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# 1 Mather theory in Lagrangian dynamics
## 1 Mather theory in Lagrangian dynamics
We recall the basics of Mather theory and state our main result, Theorem 1.10. The original references for most of the material presented in this section are Mather’s papers and . The central object is the Peierl’s barrier, introduced by Mather in . Our presentation is also influenced by the work of Fathi .
### 1.1
In this section, we consider a $`C^2`$ Hamiltonian function $`H:T^{}M\times 𝕋`$, where $`M`$ is a compact connected manifold without boundary, and $`𝕋=/`$. We denote by $`P=(q,p)`$ the points of $`T^{}M`$. The cotangent bundle is endowed with its canonical one-form $`\eta =pdq`$, and with its canonical symplectic form $`\omega =d\eta `$. Following a very standard device, we reduce our non-autonomous Hamiltonian function $`H`$ to an autonomous one by considering the extended phase space $`T^{}(M\times 𝕋)=T^{}M\times T^{}𝕋`$. We denote by $`(P,t,E)`$, $`PT^{}M`$, $`(t,E)T^{}𝕋`$ the points of this space. We consider the canonical one-form $`\lambda =pdq+Edt`$ and the associated symplectic form $`\mathrm{\Omega }=d\lambda `$. We define the new Hamiltonian $`G:T^{}(M\times 𝕋)`$ be the expression
$$G(P,t,E)=E+H(P,t).$$
We denote by $`V_G(P,t,E)`$ the Hamiltonian vector-field of $`G`$, which is defined by the relation
$$\mathrm{\Omega }_{(P,t,E)}(V_G,.)=dG_{(P,t,E)}.$$
We fix once and for all a Riemannian metric on $`M`$, and use it to define norms of tangent vectors and tangent covectors of $`M`$. We will denote this norm indifferently by $`|P|`$ or by $`|p|`$ when $`P=(q,p)T_q^{}M`$. We denote by $`\pi `$ the canonical projections $`T^{}MM`$ or $`T^{}(M\times 𝕋)M\times 𝕋`$. The theory of Mather relies on the following standard set of hypotheses.
1. Completeness. The Hamiltonian vector-field $`V_G`$ on $`T^{}(M\times 𝕋)`$ generates a complete flow, denoted by $`\mathrm{\Phi }_t`$. The flow $`\mathrm{\Phi }_t`$ preserves the level sets of $`G`$.
2. Convexity. For each $`(q,t)M\times 𝕋`$, the function $`pH(q,p,t)`$ is convex on $`T_q^{}M`$, with positive definite Hessian. Shortly, $`_p^2H>0`$.
3. Super-linearity. For each $`(q,t)M\times 𝕋`$, the function $`pH(q,p,t)`$ is super-linear, which means that $`lim_{|p|\mathrm{}}H(t,x,p)/|p|=\mathrm{}.`$
### 1.2
We associate to the Hamiltonian $`H`$ a Lagrangian function $`L:TM\times 𝕋`$ defined by
$$L(t,q,v)=\underset{pT_q^{}M}{sup}p(v)H(t,q,p).$$
The Lagrangian satisfies:
1. Convexity. For each $`(q,t)M\times 𝕋`$, the function $`vL(q,v,t)`$ is a convex function on $`T_qM`$, with positive definite Hessian. Shortly, $`_v^2L>0`$.
2. Super-linearity. For each $`(q,t)M\times 𝕋`$, the function $`vL(q,v,t)`$ is super-linear on $`T_qM`$.
Let $`X(t)=(P(t),s+t,E(t))`$ be a Hamiltonian orbit of $`G`$, and let $`q(t)=\pi (P(t))`$. Then we have the identities
$$\lambda _{X(t)}(\dot{X}(t))G(X(t))=\eta _{P(t)}(\dot{P}(t))H(P(t),s+t)=L(q(t),\dot{q}(t),s+t).$$
### 1.3
Following John Mather, we define the function $`F:M\times 𝕋\times M\times ^+`$ by
$$F(q_0,t;q_1,s)=\underset{\gamma }{\mathrm{min}}_0^sL(\gamma (\sigma ),\dot{\gamma }(\sigma ),t+\sigma )𝑑\sigma ,$$
where the minimum is taken on the set of absolutely continuous curves $`\gamma :[0,s]M`$ which satisfy $`\gamma (0)=q_0`$ and $`\gamma (1)=q_1`$. We also define the Peierl’s barrier $`h:M\times 𝕋\times M\times 𝕋\{\pm \mathrm{}\}`$ by
$$h(q_0,t_0;q_1,t_1):=\underset{n}{lim\; inf}F(q_0,t_0;q_1,s_1+n),$$
where $`t_0+s_1\text{ mod }1=t_1`$. This barrier is the central object in Mather’s study of globally minimizing orbits.
### 1.4
Let us set $`m(H)=inf_{(q,t)M\times 𝕋}h(q,t;q,t)`$. It follows from , see also , that $`m(H)\{\mathrm{},0,+\mathrm{}\}`$. In addition, for each Hamiltonian $`H`$ satisfying the hypotheses 1.1, there exists one and only one real number $`\alpha (H)`$ such that $`m(H\alpha (H))=0`$. As a consequence, there is no loss of generality in assuming that $`m(H)=0`$, or equivalently that $`\alpha (H)=0`$. We will make this assumption from now on in this section. Let us mention the terminology of Mañé, who called super-critical the Hamiltonians $`H`$ satisfying $`m(H)=+\mathrm{}`$, sub-critical the Hamiltonians satisfying $`m(H)=\mathrm{}`$, and critical the Hamiltonians satisfying $`m(H)=0`$.
### 1.5
If $`m(H)=0`$, the function $`h`$ is a real valued Lipschitz function on $`M\times 𝕋\times M\times 𝕋`$, which satisfies the triangle inequality
$$h(q_0,t_0;q_2,t_2)h(q_0,t_0;q_1,t_1)+h(q_1,t_1;q_2,t_2)$$
for all $`(q_0,t_0)`$, $`(q_1,t_1)`$ and $`(q_2,t_2)`$ in $`M\times 𝕋`$. In addition, for each $`(q,t)M\times 𝕋`$, the function $`h(q,t;.,.)`$ is a weak KAM solution in the sense of Fathi, which means that, for $`\tau \theta `$ in $``$, and $`xM`$, we have
$$h(q,t;x,\tau \text{ mod }1)=\mathrm{min}\left(h(q,t;q(\theta ),\theta \text{ mod }1)+_\theta ^\tau L(q(s),\dot{q}(s),s)𝑑s\right)$$
where the minimum is taken on the set of absolutely continuous curves $`q(s):[\theta ,\tau ]M`$ such that $`q(\tau )=x`$. Similarly, we have, for $`\tau \theta `$ in $``$, and $`xM`$,
$$h(x,\theta \text{ mod }1;q,t)=\mathrm{min}\left(h(q(\tau ),\tau \text{ mod }1;q,t)+_\theta ^\tau L(q(s),\dot{q}(s),s)𝑑s\right)$$
where the minimum is taken on the set of absolutely continuous curves $`q(s):[\theta ,\tau ]M`$ such that $`q(\theta )=x`$.
### 1.6
The projected Aubry set $`𝒜(H)`$ is the set of points $`(q,t)M\times 𝕋`$ such that $`h(q,t;q,t)=0`$. Albert Fathi proved that, for each point $`(q,t)𝒜(H)`$, the function $`h(q,t;.,.)`$ is differentiable at $`(q,t)`$. Let us denote by $`X(q,t)`$ the differential $`_3h(q,t;q,t)T_q^{}M`$ of the function $`h(q,t;.,t)`$ at point $`q`$. The Aubry set $`\stackrel{~}{}𝒜(H)`$ is defined as
$$\stackrel{~}{}𝒜(H)=\{(X(q,t),t,H(X(q,t),t));(q,t)𝒜(H)\}T^{}(M\times 𝕋).$$
The Aubry set is compact, $`\mathrm{\Phi }`$-invariant, and it is a Lipschitz graph over the projected Aubry set $`𝒜(H)`$. These are results of John Mather, see . In our presentation, which follows Fathi, this amounts to say that the function $`(q,t)X(q,t)`$ is Lipschitz on $`𝒜(H)`$.
### 1.7
The Mather set $`\stackrel{~}{}(H)`$ is defined as the union of the supports of all $`\mathrm{\Phi }`$-invariant probability measures on $`T^{}(M\times 𝕋)`$ concentrated on $`\stackrel{~}{}𝒜(H)`$. This set was first defined by Mather, but our definition is due to Mañé.
### 1.8
The projected Mañé set $`𝒩(H)`$ is the set of points $`(q,t)M\times 𝕋`$ such that there exist points $`(q_0,t_0)`$ and $`(q_1,t_1)`$ in $`𝒜(H)`$, satisfying
$$h(q_0,t_0;q_1,t_1)=h(q_0,t_0;q,t)+h(q,t;q_1,t_1).$$
Let us denote by $`(q_0,t_0;q_1,t_1)`$ the set of points $`(q,t)M\times 𝕋`$ which satisfy this relation. If $`(q_0,t_0)𝒜(H)`$ and $`(q_1,t_1)𝒜(H)`$ are given, and if $`(q,t)(q_0,t_0;q_1,t_1)`$, then the function $`h(q_0,t_0;.,t)`$ is differentiable at $`q`$, as well as the function $`h(.,t;q_1,t_1)`$, and $`_3h(q_0,t_0,q,t)+_1h(q,t;q_1,t_1)=0`$. This is proved in following ideas of Albert Fathi. We define
$$\stackrel{~}{}(q_0,t_0;q_1,t_1):=\left\{(_3h(q_0,t_0,q,t),t,H(_3h(q_0,t_0,q,t),t)),(q,t)(q_0,t_0;q_1,t_1)\right\}.$$
The set $`\stackrel{~}{}(q_0,t_0;q_1,t_1)`$ is a compact $`\mathrm{\Phi }`$-invariant subset of $`T^{}(M\times 𝕋)`$, and it is a Lipschitz Graph. The Mañé set $`\stackrel{~}{}𝒩(H)`$ is the set
$$\stackrel{~}{}𝒩(H)=\underset{(q_0,t_0),(q_1,t_1)𝒜(H)}{}\stackrel{~}{}(q_0,t_0;q_1,t_1)T^{}(M\times 𝕋).$$
The Mañé set was first introduced by Mather in , it is compact and $`\mathrm{\Phi }`$-invariant, and it contains the Aubry set. In other words, we have the important inclusions
$$\stackrel{~}{}(H)\stackrel{~}{}𝒜(H)\stackrel{~}{}𝒩(H).$$
The Mañé set is usually not a graph. However, it satisfies
$$\stackrel{~}{}𝒩(H)\pi ^1\left(𝒜(H)\right)=\stackrel{~}{}𝒜(H).$$
This follows from the fact, proved by Albert Fathi, that, for each $`(x,\theta )M\times 𝕋`$ and each $`(q,t)𝒜(H)`$, the function $`h(x,\theta ;.,t)`$ is differentiable at $`q`$ and satisfies $`_3h(x,\theta ;q,t)=X(q,t)`$.
### 1.9
Mather introduced the function $`d(q,t;q^{},t^{})=h(q,t;q^{},t^{})+h(q^{},t^{};q,t)`$ on $`M\times 𝕋`$. When restricted to $`𝒜(H)\times 𝒜(H)`$, it is a pseudo-metric. This means that this function is symmetric, non-negative, satisfies the triangle inequality, and $`d(q,t;q,t)=0`$ for $`(q,t)𝒜(H)`$. We shall also denote by $`d`$ the pseudo-metric $`d(P,t,H(P,t);P^{},t^{},H(P^{},t^{}))=d(\pi (P),t;\pi (P^{}),t^{})`$ on $`\stackrel{~}{A}(H)`$. The relation $`d(P,t,E;P^{},t^{},E^{})=0`$ is an equivalence relation on $`\stackrel{~}{}𝒜(H)`$. The classes of equivalence are called the static classes. Let us denote by $`\dot{}𝒜(H)`$ the set of static classes. The pseudo-metric $`d`$ gives rise to a metric $`\dot{d}`$ on $`\dot{}𝒜(H)`$. The compact metric space $`(\dot{}𝒜(H),\dot{d})`$ is called the quotient Aubry set. It was introduced by John Mather.
### 1.10
The diffeomorphism $`\mathrm{\Psi }:T^{}(M\times 𝕋)T^{}(M\times 𝕋)`$ is called exact if the form $`\mathrm{\Psi }^{}\lambda \lambda `$ is exact.
Theorem Let $`H`$ be a Hamiltonian satisfying the hypotheses 1.1, and let $`\mathrm{\Psi }:T^{}(M\times 𝕋)T^{}(M\times 𝕋)`$ be an exact diffeomorphism such that the Hamiltonian
$$\mathrm{\Psi }^{}H:=G\mathrm{\Psi }(P,t,E)E$$
is independent of $`E`$ and satisfies the hypotheses 1.1 when considered as a function on $`T^{}M\times 𝕋`$. Then $`m(\mathrm{\Psi }^{}H)=m(H)`$ hence $`\alpha (H)=\alpha (\mathrm{\Psi }^{}H)`$. If $`m(H)=0`$, then we have
$$\mathrm{\Psi }(\stackrel{~}{}(\mathrm{\Psi }^{}H))=\stackrel{~}{}(H),\mathrm{\Psi }(\stackrel{~}{}𝒜(\mathrm{\Psi }^{}H))=\stackrel{~}{}𝒜(H),\mathrm{\Psi }(\stackrel{~}{}𝒩(\mathrm{\Psi }^{}H))=\stackrel{~}{}𝒩(H).$$
In addition, $`\mathrm{\Psi }`$ sends the static classes of $`\mathrm{\Psi }^{}H`$ onto the static classes of $`H`$, and the induced mapping
$$\dot{\mathrm{\Psi }}:\dot{}𝒜(\mathrm{\Psi }^{}H)\dot{}𝒜(H)$$
is an isometry for the quotient metrics.
### 1.11
We prove this result in the sequel. In section 2, we set the basis of a symplectic Aubry-Mather theory for general Hamiltonian systems. We prove that the analogue of Theorem 1.10 holds in this general setting. We also continue the theory a bit further than would be necessary to prove Theorem 1.10. In section 3, we prove that, under the hypotheses of Theorem 1.10, the symplectic Aubry-Mather sets coincide with the standard Aubry-Mather sets, which ends the proof of Theorem 1.10.
## 2 A barrier in phase space
We propose general definitions for a Mather theory of Hamiltonian systems. Of course, the definitions given below provide relevant objects only for some specific Hamiltonian systems. It would certainly be interesting to give natural conditions on $`H`$ implying non-triviality of the theory developed in this section. We shall only check, in the next section, that our definitions coincide with the standard ones in the convex case, obtaining non-triviality in this special case. Let us mention once again that it might be possible and interesting to find more geometric definition using the methods of .
### 2.1
In this section, we work in a very general setting. We consider a manifold $`N`$, not necessarily compact, and an autonomous Hamiltonian function $`G:T^{}N`$. We assume that $`G`$ generates a complete Hamiltonian flow $`\mathrm{\Phi }_t`$. We make no convexity assumption. We denote by $`\lambda `$ the canonical one-form of $`T^{}N`$, and by $`V_G(P)`$ the Hamiltonian vector-field of $`G`$. Let $`D(P,P^{})`$ be a distance on $`T^{}N`$ induced from a Riemannian metric. We identify $`N`$ with the zero section of $`T^{}N`$, so that $`D`$ is also a distance on $`N`$. We assume that $`D(\pi (X),\pi (X^{}))D(X,X^{})`$ for $`X`$ and $`X^{}`$ in $`T^{}N`$.
### 2.2
Let $`X_0`$ and $`X_1`$ be two points of $`T^{}N`$. A pre-orbit between $`X_0`$ and $`X_1`$ is the data of a sequence $`\underset{¯}{Y}=(Y_n)`$ of curves $`Y_n(s):[0,T_n]T^{}N`$ such that:
1. For each $`n`$, the curve $`Y_n`$ has a finite number $`N_n`$ of discontinuity points $`T_n^i]0,T_n[,1iN_n`$ such that $`T_n^{i+1}>T_n^i`$. We shall also often use the notations $`T_n^0=0`$ and $`T_n^{N_n+1}=T_n`$.
2. The curve $`Y_n`$ satisfies $`Y_n(T_n^i+s)=\mathrm{\Phi }_s(Y_n(T_n^i))`$ for each $`s[0,T_n^{i+1}T_n^i[`$. We denote by $`Y_n(T_n^i)`$ the point $`\mathrm{\Phi }_{T_n^iT_n^{i1}}(Y(T_n^{i1}))`$ and impose that $`Y_n(T_n)=Y_n(T_n)`$.
3. We have $`T_n\mathrm{}`$ as $`n\mathrm{}`$.
4. We have $`Y_n(0)X_0`$ and $`Y_n(T_n)X_1`$. In addition, we have $`lim_n\mathrm{}\mathrm{\Delta }(Y_n)=0`$, where we denote by $`\mathrm{\Delta }(Y_n)`$ the sum $`_{i=1}^{N_n}D(Y_n(T_n^i),Y_n(T_n^i))`$.
5. There exists a compact subset $`KT^{}N`$ which contains the images of all the curves $`Y_n`$.
The pre-orbits do not depend on the metric which has been used to define the distance $`D`$. In a standard way, we call action of the curve $`Y_n(t)`$ the value
$$A(Y_n)=_0^{T_n}\lambda _{Y_n(t)}(\dot{Y}_n(t))G(Y_n(t))dt.$$
The action of the pre-orbit $`\underset{¯}{Y}`$ is
$$A(\underset{¯}{Y}):=\underset{n\mathrm{}}{lim\; inf}A(Y_n).$$
### 2.3
Lemma If there exists a pre-orbit between $`X_0`$ and $`X_1`$, then $`G(X_0)=G(X_1)`$.
Proof. This follows easily from the fact that the Hamiltonian flow $`\mathrm{\Phi }`$ preserves the Hamiltonian function $`G`$.
### 2.4
We define the barrier $`\stackrel{~}{h}:T^{}M\times T^{}M\{\pm \mathrm{}\}`$ by the expression
$$\stackrel{~}{h}(X_0,X_1)=\underset{\underset{¯}{Y}}{inf}A(\underset{¯}{Y})$$
where the infimum is taken on the set of pre-orbits between $`X_0`$ and $`X_1`$. As usual, we set $`\stackrel{~}{h}(X_0,X_1)=+\mathrm{}`$ if there does not exist any pre-orbit between $`X_0`$ and $`X_1`$. If $`\stackrel{~}{h}(X_0,X_1)<+\mathrm{}`$, then the forward orbit of $`X_0`$ and the backward orbit of $`X_1`$ are bounded. As a consequence, if $`\stackrel{~}{h}(X,X)<+\mathrm{}`$, then the orbit of $`X`$ is bounded.
### 2.5
Property For each $`t>0`$, we have the equality
$$\stackrel{~}{h}(X_0,X_1)=\stackrel{~}{h}(\mathrm{\Phi }_t(X_0),X_1)+_0^t\lambda _{\mathrm{\Phi }_s(X_0)}\left(V_G(\mathrm{\Phi }_s(X_0))\right)G(\mathrm{\Phi }_s(X_0))ds$$
and
$$\stackrel{~}{h}(X_0,\mathrm{\Phi }_t(X_1))=\stackrel{~}{h}(X_0,X_1)+_0^t\lambda _{\mathrm{\Phi }_s(X_1)}\left(V_G(\mathrm{\Phi }_s(X_1))\right)G(\mathrm{\Phi }_s(X_1))ds$$
Proof. We shall prove the first equality, the proof of the second one is similar. To each pre-orbit $`\underset{¯}{Y}`$ between $`X_0`$ and $`X_1`$, we associate the pre-orbit $`\underset{¯}{Z}`$ between $`\mathrm{\Phi }_t(X_0)`$ and $`X_1`$ defined by $`Z_n(s):[0,T_nt]Y_n(s+t)`$. We have
$$A(\underset{¯}{Y})=A(\underset{¯}{Z})+_0^t\lambda _{\mathrm{\Phi }_s(X_0)}\left(V_G(\mathrm{\Phi }_s(X_0))\right)G(\mathrm{\Phi }_s(X_0))ds$$
This implies that
$$\stackrel{~}{h}(\mathrm{\Phi }_t(X_0),X_1)\stackrel{~}{h}(X_0,X_1)_0^t\lambda _{\mathrm{\Phi }_s(X_0)}\left(V_G(\mathrm{\Phi }_s(X_0))\right)G(\mathrm{\Phi }_s(X_0))ds.$$
In a similar way, we associate to each pre-orbit $`\underset{¯}{Z}=Z_n(s):[0,T_n]T^{}M`$ between $`\mathrm{\Phi }_t(X_0)`$ and $`X_1`$ the pre-orbits $`\underset{¯}{Y}:[0,T_n+t]T^{}M`$ between $`X_0`$ and $`X_1`$ defined by $`Y_n(s)=\mathrm{\Phi }_{st}(Z_n(0))`$ for $`s[0,t]`$ and $`Y_n(s)=Z_n(st)`$ for $`s[t,T_n+t]`$. We have
$$A(\underset{¯}{Y})=A(\underset{¯}{Z})+_0^t\lambda _{\mathrm{\Phi }_s(X_0)}\left(V_G(\mathrm{\Phi }_s(X_0))\right)G(\mathrm{\Phi }_s(X_0))ds.$$
This implies that
$$\stackrel{~}{h}(X_0,X_1)\stackrel{~}{h}(\mathrm{\Phi }_t(X_0),X_1)+_0^t\lambda _{\mathrm{\Phi }_s(X_0)}\left(V_G(\mathrm{\Phi }_s(X_0))\right)G(\mathrm{\Phi }_s(X_0))ds.$$
### 2.6
property The function $`\stackrel{~}{h}`$ satisfies the triangle inequality. More precisely, the relation
$$\stackrel{~}{h}(X_1,X_3)\stackrel{~}{h}(X_1,X_2)+\stackrel{~}{h}(X_2,X_3)$$
holds for each points $`X_1`$, $`X_2`$ and $`X_3`$ such that the right hand side has a meaning.
Proof. If one of the values $`\stackrel{~}{h}(X_1,X_2)`$ or $`\stackrel{~}{h}(X_2,X_3)`$ is $`+\mathrm{}`$, then there is nothing to prove. If they are both different from $`+\mathrm{}`$, then, for each $`ϵ>0`$ there exists a pre-orbits $`\underset{¯}{Y}=Y_n:[0,T_n]T^{}N`$ between $`X_1`$ and $`X_2`$ such that $`A(\underset{¯}{Y})\stackrel{~}{h}(X_1,X_2)+ϵ`$ (resp. $`A(\underset{¯}{Y})1/ϵ`$ in the case where $`\stackrel{~}{h}(X_1,X_2)=\mathrm{}`$) and a pre-orbits $`\underset{¯}{Y}^{}=Y_n^{}:[0,S_n]T^{}N`$ between $`X_2`$ and $`X_3`$ such that $`A(\underset{¯}{Y}^{})\stackrel{~}{h}(X_2,X_3)+ϵ`$ (resp. $`A(\underset{¯}{Y}^{})1/ϵ`$ in the case where $`\stackrel{~}{h}(X_1,X_2)=\mathrm{}`$). Let us consider the sequence of curves $`Z_n(t):[0,T_n+S_n]T^{}N`$ such that $`Z_n=X_n`$ on $`[0,T_n[`$ and $`Z_n(t+T_n)=Y_n(t)`$ for $`t[0,S_n]`$. It is clear that the sequence $`\underset{¯}{Z}=Z_n`$ is a pre-orbit between $`X_1`$ and $`X_3`$, and that its action satisfies
$$A(\underset{¯}{Z})=A(\underset{¯}{X})+A(\underset{¯}{Y})\stackrel{~}{h}(X_1,X_2)+\stackrel{~}{h}(X_2,X_3)+2ϵ.$$
As a consequence, for all $`ϵ>0`$, we have $`\stackrel{~}{h}(X_1,X_3)\stackrel{~}{h}(X_1,X_2)+\stackrel{~}{h}(X_2,X_3)+2ϵ`$ hence the triangle inequality holds.
### 2.7
property Let $`\mathrm{\Psi }:T^{}NT^{}N`$ be an exact diffeomorphism. We have the equality
$$\stackrel{~}{h}_{G\mathrm{\Psi }}(X_0,X_1)=\stackrel{~}{h}_G(\mathrm{\Psi }(X_0),\mathrm{\Psi }(X_1))+S(X_0)S(X_1),$$
where $`S:T^{}N`$ is a function such that $`\mathrm{\Psi }^{}\lambda \lambda =dS`$.
Proof. Observe first that $`\underset{¯}{Y}=Y_n`$ is a pre-orbit for the Hamiltonian $`G\mathrm{\Psi }`$ between points $`X_0`$ and $`X_1`$ if and only if $`\mathrm{\Psi }(\underset{¯}{Y})=\mathrm{\Psi }(Y_n)`$ is a pre-orbit for the Hamiltonian $`G`$ between $`\mathrm{\Psi }(X_0)`$ and $`\mathrm{\Psi }(X_1)`$. As a consequence, it is enough to prove that
$$A_{G\mathrm{\Psi }}(\underset{¯}{Y})=A_G(\mathrm{\Psi }(\underset{¯}{Y}))+S(X_0)S(X_1).$$
Let us denote by $`\underset{¯}{Z}=Z_n`$ the pre-orbit $`\mathrm{\Psi }(Y_n)`$. Setting $`T_n^0=0`$ and $`T_n^{N_n+1}=T_n`$, we have
$$A_G(Z_n)=\underset{i=0}{\overset{N_n}{}}_{T_n^i}^{T_n^{i+1}}\lambda _{Z_n(t)}(\dot{Z}_n(t))G(Z_n(t))dt$$
$$=\underset{i=0}{\overset{N_n}{}}_{T_n^i}^{T_n^{i+1}}(\mathrm{\Psi }^{}\lambda )_{Y_n(t)}(\dot{Y}_n(t))G\mathrm{\Psi }(Y_n(t))dt$$
$$=\underset{i=0}{\overset{N_n}{}}\left(_{T_n^i}^{T_n^{i+1}}\lambda _{Y_n(t)}(\dot{Y}_n(t))G\mathrm{\Psi }(Y_n(t))dt+S(Y_n(T_n^{i+1}))S(Y_n(T_n^i))\right)$$
$$=A_{G\mathrm{\Psi }}(Y_n)S(Y_n(0))+S(Y_n(T_n))+\underset{i=1}{\overset{N_n}{}}(S(Y_n(T_n^i))S(Y_n(T_n^i)).)$$
Since the function $`S`$ is Lipschitz on the compact set $`K`$ which contains the image of the curves $`Y_n`$, we obtain at the limit
$$A_G(\underset{¯}{Z})=A_{G\mathrm{\Psi }}(\underset{¯}{Y})S(X_0)+S(X_1).$$
### 2.8
proposition Let us set $`\stackrel{~}{m}(H):=inf_{XT^{}N}\stackrel{~}{h}(X,X)`$. We have $`\stackrel{~}{m}(H)\{\mathrm{},0,+\mathrm{}\}`$. In addition, if $`\stackrel{~}{m}(H)=0`$, then there exists a point $`X`$ in $`T^{}N`$ such that $`\stackrel{~}{h}(X,X)=0`$.
Proof. It follows from the triangle inequality that, for each $`XT^{}N`$, $`\stackrel{~}{h}(X,X)0`$ or $`\stackrel{~}{h}(X,X)=\mathrm{}`$. As a consequence, $`\stackrel{~}{m}(H)0`$ or $`\stackrel{~}{m}(H)=\mathrm{}`$. Let us assume that $`\stackrel{~}{m}(H)[0,\mathrm{}[`$. Then there exists a point $`X_0T^{}N`$ and a pre-orbits $`\underset{¯}{Y}=Y_n:[0,T_n]T^{}N`$ between $`X_0`$ and $`X_0`$ such that $`A(\underset{¯}{Y})[0,\mathrm{}[`$. Let $`K`$ be a compact subset of $`T^{}N`$ which contains the image of all the curves $`Y_n`$. Let $`S_n`$ be a sequence of integers such that $`T_n/S_n\mathrm{}`$ and $`S_n\mathrm{}`$. Let $`b_n`$ be the integer part of $`T_n/S_n`$. Note that $`b_n\mathrm{}`$. Let $`d_n`$ be a sequence of integers such that $`d_n\mathrm{}`$ and $`d_n/b_n0`$. Since the set $`K`$ is compact, there exists a sequence $`ϵ_n0`$ such that, whenether $`b_n`$ points are given in $`K`$, then at least $`d_n`$ of them lie in a same ball of radius $`ϵ_n`$. So there exists a point $`X_nK`$ such that at least $`d_n`$ of the points $`Y_n(S_n),Y_n(2S_n),\mathrm{},Y_n(b_nS_n)`$ lie in the ball of radius $`ϵ_n`$ and center $`X_n`$. Let us denote by $`Y_n(t_n^1),Y_n(t_n^2)\mathrm{},Y_n(t_n^{d_n})`$ these points, where $`t_n^{i+1}t_n^i+S_n`$. We can assume, taking a subsequence, that the sequence $`X_n`$ has a limit $`X`$ in $`K`$. It is not hard to see that $`\underset{¯}{Y}^i=Y_{n|[t_n^i,t_n^{i+1}]}`$ is a pre-orbit between $`X`$ and $`X`$. On the other hand, for each $`k`$, we define the sequence of curves $`Z_n^k:[0,T_n+t_n^1t_n^k]T^{}N`$ by $`Z_n^k(t)=Y_n(t)`$ for $`t[0,t_n^1[`$, and $`Z_n^k(t)=Y_n(t+t_n^kt_n^1)`$ for $`t[t_n^1,T_n+t_n^1t_n^k]`$. For each $`k`$, the sequence $`Z_n^k`$ is a pre-orbit between $`X_0`$ and $`X_0`$. We have
$$A(Y_n)=A(Z_n^k)+\underset{i=1}{\overset{k1}{}}A(Y_n^i)$$
hence
$$A(\underset{¯}{Y})\stackrel{~}{h}(X_0,X_0)+(k1)\stackrel{~}{h}(X,X).$$
Since $`A(\underset{¯}{Y})`$ is a real number, and since this inequality holds for all $`k`$, this implies that $`\stackrel{~}{h}(X,X)=0`$.
### 2.9
Let us define the symplectic Aubry set of $`G`$ as the set
$$\stackrel{~}{}𝒜_s(G):=\{XT^{}N\text{ such that }\stackrel{~}{h}(X,X)=0\text{ and }G(X)=0\}T^{}N.$$
The symplectic Mather set $`\stackrel{~}{}_s(G)`$ of $`G`$ is the union of the supports of the compactly supported $`\mathrm{\Phi }`$-invariant probability measures concentrated on $`\stackrel{~}{}𝒜_s(G)`$. Note that, in general, it is not clear that the symplectic Aubry set should be closed. The symplectic Mather set, then, may not be contained in the symplectic Aubry set, but only in its closure. The Mather set and the Aubry set are $`\mathrm{\Phi }`$-invariant, as follows directly from 2.5. If $`\stackrel{~}{m}(H)=0`$, then the symplectic Aubry set is not empty, and all its orbits are bounded, hence the symplectic Mather set $`\stackrel{~}{}_s(G)`$ is not empty.
### 2.10
For each pair $`X_0`$, $`X_1`$ of points in $`\stackrel{~}{}𝒜_s(G)`$, we define the set $`\stackrel{~}{}_s(X_0,X_1)`$ of points $`PT^{}N`$ such that
$$\stackrel{~}{h}(X_0,X_1)=\stackrel{~}{h}(X_0,X)+\stackrel{~}{h}(X,X_1)$$
if $`\stackrel{~}{h}(X_0,X_1)`$, and $`\stackrel{~}{}_s(X_0,X_1)=\mathrm{}`$ otherwise. Note that the sets $`\stackrel{~}{}_s(X_0,X_1)`$ are all contained in the level $`\{G=0\}`$. Indeed, the finiteness of $`\stackrel{~}{h}(X_0,X)`$ implies that $`G(X_0)=G(X)`$, while $`G(X_0)=0`$ by definition of $`\stackrel{~}{}𝒜_s(G)`$. It follows from 2.5 that the set $`\stackrel{~}{}_s(X_0,X_1)`$ is $`\mathrm{\Phi }`$-invariant. We now define the symplectic Mañé set as
$$\stackrel{~}{}𝒩_s(G):=\underset{X_0,X_1\stackrel{~}{}𝒜_s(G)}{}\stackrel{~}{}_s(X_0,X_1).$$
The Mañé set is $`\mathrm{\Phi }`$-invariant, all its orbits are bounded. We have the inclusion
$$\stackrel{~}{}𝒜_s(G)\stackrel{~}{}𝒩_s(G).$$
In order to prove this inclusion, just observe that $`X_0\stackrel{~}{}(X_0,X_0)`$ for each $`X_0\stackrel{~}{}𝒜_s(G)`$.
### 2.11
If $`\mathrm{\Psi }:T^{}NT^{}N`$ is an exact diffeomorphism, then we have
$$\mathrm{\Psi }(\stackrel{~}{}_s(G\mathrm{\Psi }))=\stackrel{~}{}_s(G),\mathrm{\Psi }(\stackrel{~}{}𝒜_s(G\mathrm{\Psi }))=\stackrel{~}{}𝒜_s(G),\mathrm{\Psi }(\stackrel{~}{}𝒩_s(G\mathrm{\Psi }))=\stackrel{~}{}𝒩_s(G),$$
this follows obviously from 2.7, and from the fact that $`\mathrm{\Psi }`$ conjugates the Hamiltonian flow of $`G`$ and the Hamiltonian flow of $`G\mathrm{\Psi }`$.
### 2.12
Let us assume that $`\stackrel{~}{m}(G)=0`$, and set
$$\stackrel{~}{d}(X,X^{})=\stackrel{~}{h}(X,X^{})+\stackrel{~}{h}(X^{},X).$$
We have $`\stackrel{~}{d}(X,X^{})0`$, and the function $`\stackrel{~}{d}`$ satisfies the triangle inequality, and is symmetric. In addition, we obviously have $`\stackrel{~}{d}(X,X)=0`$ if and only if $`X\stackrel{~}{}𝒜_s(G)`$. The restriction of the function $`\stackrel{~}{d}`$ to the set $`\stackrel{~}{}𝒜_s(G)`$ is a pseudo-metric with $`+\mathrm{}`$ as a possible value. We define an equivalence relation on $`\stackrel{~}{}𝒜_s(G)`$ by saying that the points $`X`$ and $`X^{}`$ are equivalent if and only if $`\stackrel{~}{d}(X,X^{})=0`$. The equivalence classes of this relation are called the static classes. Let us denote by $`(\dot{}𝒜_s(G),\dot{d}_s`$) the metric space obtained from $`\stackrel{~}{}𝒜_s`$ by identifying points $`X`$ and $`X^{}`$ when $`\stackrel{~}{d}(X,X^{})=0`$. In other words, the set $`\dot{}𝒜_s(G)`$ is the set of static classes of $`H`$. We call $`(\dot{}𝒜_s(G),\dot{d}_s)`$ the quotient Aubry set. Note that the metric $`\dot{d}_s`$ can take the value $`+\mathrm{}`$. The quotient Aubry set is also well behaved under exact diffeomorphisms. More precisely, if $`\mathrm{\Psi }`$ is an exact diffeomorphism of $`T^{}N`$, then the image of a static class of $`G\mathrm{\Psi }`$ is a static class of $`G`$. This defines a map
$$\dot{\mathrm{\Psi }}:\dot{}𝒜_s(G\mathrm{\Psi })\dot{}𝒜_s(G)$$
which is an isometry for the quotient metrics.
### 2.13
Proposition, Assume that $`\stackrel{~}{m}(G)=0`$, and in addition that the function $`\stackrel{~}{h}`$ is bounded from below. Then the orbits of $`\stackrel{~}{}𝒩_s(G)`$ are bi-asymptotic to $`\stackrel{~}{}𝒜_s(G)`$. In addition, for each orbit $`X(s)`$ in $`\stackrel{~}{}𝒩_s(G)`$, there exists a static class $`S`$ in $`\stackrel{~}{}𝒜_s(G)`$ and a static class $`S+`$ such that the orbit $`X(s)`$ is $`\alpha `$-asymptotic to $`S`$ and $`\omega `$-asymptotic to $`S+`$.
Proof. Let $`\omega `$ and $`\omega ^{}`$ be two points in the $`\omega `$-limit of the orbits $`X(t)=\mathrm{\Phi }_t(X)`$. We have to prove that $`\omega `$ and $`\omega ^{}`$ belong to the symplectic Aubry set, and to the same static class. It is enough to prove that $`\stackrel{~}{d}(\omega ,\omega ^{})=0`$. In order to do so, we consider two increasing sequences $`t_n`$ and $`s_n`$, such that $`t_ns_n\mathrm{}`$, $`s_nt_{n1}\mathrm{}`$, $`X(t_n)\omega `$ and $`X(s_n)\omega ^{}`$. Let $`\underset{¯}{Y}=Y_n:[0,t_ns_n]T^{}N`$ be the pre-orbit between $`\omega ^{}`$ and $`\omega `$ defined by $`Y_n(t)=X(ts_n)`$. Similarly, we consider the pre-orbit $`\underset{¯}{Z}=Z_n:[0,s_{n+1}t_n]T^{}N`$ between $`\omega `$ and $`\omega ^{}`$ defined by $`Z_n(t)=X(tt_n)`$. Since $`X`$ belongs to $`\stackrel{~}{}𝒩_s(G)`$, there exist points $`X_0`$ and $`X_1`$ in $`\stackrel{~}{}𝒜_s(G)`$ such that $`X\stackrel{~}{}(X_0,X_1)`$. In view of 2.5, we have
$$\stackrel{~}{h}(X(t_n),X_1)=\stackrel{~}{h}(X(t_m),X_1)+_{t_n}^{t_m}\lambda _{X(t)}(\dot{X}(t))G(X(t))dt$$
for all $`mn`$. Since the function $`\stackrel{~}{h}`$ is bounded from below, we conclude that the double sequence $`_{t_n}^{t_m}\lambda _{X(t)}(\dot{X}(t))G(X(t))dt,mn`$ is bounded from above, so that
$$lim\; inf_{t_n}^{t_{n+1}}\lambda _{X(t)}(\dot{X}(t))G(X(t))dt0.$$
As a consequence, we have $`lim\; infA(Y_{n+1})+A(Z_n)0`$ hence $`A(\underset{¯}{Y})+A(\underset{¯}{Z})=0`$, and $`\stackrel{~}{d}(\omega ,\omega ^{})=0`$. The proof is similar for the $`\alpha `$-limit.
It is useful to finish with section with a technical remark.
### 2.14
Lemma Let $`\underset{¯}{Y}=Y_n:[0,T_n]T^{}N`$ be a pre-orbit between between $`X_0`$ and $`X_1`$. There exists a pre-orbit $`\underset{¯}{Z}`$ between $`X_0`$ and $`X_1`$ which has the same action as $`\underset{¯}{Y}`$, and has discontinuities only at times $`1,2,\mathrm{},[T_n]1`$, where $`[T_n]`$ is the integer part of $`T_n`$.
Proof. We set $`Z_n(k+s)=\mathrm{\Phi }_s(Y_n(k))`$ for each $`k=0,1,\mathrm{},[T_n]2`$, and $`s[0,1[`$, and $`Z_n([T_n]1+s)=\mathrm{\Phi }_s(Y_n([T_n]1))`$ for each $`s[0,1+T_n[T_n][`$. It is not hard to see that $`A(Z_n)A(Y_n)0`$, hence $`A(\underset{¯}{Y})=A(\underset{¯}{Z})`$.
## 3 The case of convex Hamiltonian systems
We assume the hypotheses 1.1, and prove that the symplectic definitions of section 2 agree with the standard definitions of section 1. This proves that the theory of section 2 is not trivial at least in this case. This also ends the proof of Theorem 1.10.
### 3.1
In this section, we consider a Hamiltonian function $`H:T^{}M\times 𝕋`$ satisfying the hypotheses 1.1. We set $`N=M\times 𝕋`$. We denote by $`(P,t,E)`$ the points of $`T^{}N`$ and set $`G(P,t,E)=E+H(P,t):T^{}N`$. We denote by $`h(q,t;q^{},t^{})`$ the Peierl’s barrier associated to $`H`$ in section 1 and by $`\stackrel{~}{h}(P,t,E;P^{},t^{},E^{})`$ the barrier associated to $`G`$ in section 2.
### 3.2
Before we state the main result of this section, some terminology is necessary. If $`u:M`$ is a continuous function, we say that $`PT_q^{}M`$ is a proximal super-differential of $`u`$ at point $`q`$ (or simply a super-differential) if there exists a smooth function $`f:M`$ such that $`fu`$ has a minimum at $`q`$ and $`df_q=P`$. Clearly, if $`u`$ is differentiable at $`q`$ and if $`P`$ is a proximal super-differential of $`u`$ at $`q`$, then $`P=du_q`$.
### 3.3
Proposition We have the relation
$$h(q,t;q^{},t^{})=\underset{PT_q^{}M,P^{}T_q^{}^{}M}{\mathrm{min}}\stackrel{~}{h}(P,t,H(P,t);P^{},t^{},H(P^{},t^{})).$$
In addition, if the minimum is reached at $`(P,P^{})`$ then $`P`$ is a super-differential of the function $`h(.,t;q^{},t^{})`$ at point $`q`$ and $`P^{}`$ is a super-differential of the function $`h(q,t;.,t^{})`$ at point $`q^{}`$.
Proof. Let us fix two points $`(q,t)`$ and $`(q^{},t^{})`$ in $`N=M\times 𝕋`$. We claim that the inequality
$$\stackrel{~}{h}(P,t,E;P^{},t^{},E^{})h(q,t;q^{},t^{})$$
holds for each $`(P,t,E)T_{(q,t)}^{}N`$ and each $`(P^{},t^{},E^{})T_{(q^{},t^{})}^{}N`$. If $`\stackrel{~}{h}(P,t,E;P^{},t^{},E^{})=+\mathrm{}`$, then there is nothing to prove. Else, let us fix $`ϵ>0`$. There exists a pre-orbit $`\underset{¯}{Y}=Y_n(s):[0,T_n]T^{}N`$ between $`(P,t,E)`$ and $`(P^{},t^{},E^{})`$ such that $`A(\underset{¯}{Y})\stackrel{~}{h}(P,t,E;P^{},t^{},E^{})+ϵ`$ (resp. $`A(\underset{¯}{Y})1/ϵ`$ in the case where $`\stackrel{~}{h}(P,t,E;P^{},t^{},E^{})=\mathrm{}`$). In view of 2.14, it is possible to assume that the discontinuity points $`T_n^i`$ of $`Y_n`$ satisfy $`T_n^{i+1}T_n^i+1`$. Let us write
$$Y_n(s)=(P_n(s),\tau _n(s),E_n(s)),$$
and $`q_n(s)=\pi (P_n(s))`$. Let $`\delta _n^i`$ be the real number closest to $`T_n^{i+1}T_n^i`$ among those which satisfy $`\tau _n(T_n^i)+\delta _n^i=\tau _n(T_n^{i+1})`$.
We have
$$A(Y_n)=\underset{i=0}{\overset{N_n}{}}_{T_n^i}^{T_n^{i+1}}L(q_n(s),\dot{q}_n(s),s+\tau _n(T_n^i))𝑑t\underset{i=0}{\overset{N_n}{}}F(q(T_n^i),\tau _n(T_n^i);q(T_n^{i+1}),T_n^{i+1}T_n^i).$$
It is known that the functions $`F(q,t;q^{},s)`$ is Lipschitz on $`\{s1\}`$, see for example , 3.2. We have
$$\underset{i=0}{\overset{N_n}{}}\left|F(q_n(T_n^i),\tau _n(T_n^i);q_n(T_n^{i+1}),T_n^{i+1}T_n^i)F(q_n(T_n^i),\tau _n(T_n^i);q_n(T_n^{i+1}),\delta _n^i)\right|$$
$$C\underset{i=0}{\overset{N_n1}{}}D(q_n(T_n^{i+1}),\tau _n(T_n^{i+1});q_n(T_n^{i+1}),\tau _n(T_n^{i+1}))$$
$$C\underset{i=0}{\overset{N_n1}{}}D(Y_n(T_n^{i+1}),Y_n(T_n^{i+1}))0.$$
As a consequence, we have
$$A(\underset{¯}{Y})lim\; inf\underset{i=0}{\overset{N_n}{}}F(q(T_n^i),\tau _n(T_n^i);q(T_n^{i+1}),\delta _n^i)$$
$$lim\; infF(q_n(0),\tau _n(0);q_n(T_n),\underset{i=0}{\overset{N_n}{}}\delta _n^i)h(q,t;q^{},t^{}),$$
hence $`ϵ+\stackrel{~}{h}(P,t,E;P^{},t^{},E^{})h(q,t;q^{},t^{})`$ (resp. $`1/ϵh(q,t;q^{},t^{})`$). Since this holds for all $`ϵ>0`$, we have $`\stackrel{~}{h}(P,t,E;P^{},t^{},E^{})h(q,t;q^{},t^{})`$ as desired.
Conversely, let us consider a sequence $`T_n`$ such that $`T_n\mathrm{}`$, $`t+T_n\text{ mod }1=t^{}`$, and
$$h(q,t;q^{},t^{})=\underset{n\mathrm{}}{lim}F(q,t;q^{},T_n).$$
Let $`q_n(s):[0,T_n]M`$ be a curve such that
$$_0^{T_n}L(q_n(s),\dot{q}_n(s),s+t)𝑑s=F(q,t;q^{},T_n).$$
Since the curve $`q_n`$ is minimizing the action, there exists a Hamiltonian trajectory
$$Y_n(s)=(P_n(s),t+s,E_n(s)=H(X_n(s),t+s)):[0,T_n]T^{}N$$
whose projection on $`M`$ is the curve $`q_n`$. In addition, by well known results on minimizing orbits, see , there exists a compact subset of $`T^{}M`$ which contains the images of all the curves $`P_n(s)`$. As a consequence, we can assume, taking a subsequence if necessary, that the sequences $`P_n(0)`$ and $`P_n(T_n)`$ have limits $`PT_q^{}M`$ and $`P^{}T_q^{}^{}M`$. The sequence $`\underset{¯}{Y}=Y_n`$ is then a pre-orbit between $`(P,t,H(P,t))`$ and $`(P^{},t^{},H(P^{},t^{}))`$, and its action is
$$A(\underset{¯}{Y})=limA(Y_n)=lim_0^{T_n}L(q_n(s),\dot{q}_n(s),t+s)𝑑s=h(q,t;q^{},t^{}).$$
As a consequence, we have
$$\stackrel{~}{h}(P,t,H(P,t));P^{},t^{},H(P^{},t^{}))h(q,t;q^{},t^{}).$$
This ends the proof of the first part of the Proposition.
Let now $`Y=(P,t,E)T_q^{}M\times T^{}𝕋`$ and $`Y^{}=(P^{},t^{},E^{})T_q^{}^{}M\times T^{}𝕋`$ be points such that $`h(q,t;q^{},t^{})=\stackrel{~}{h}(Y;Y^{})`$. Let $`q(s)`$ be the projection on $`M`$ of the orbit $`\mathrm{\Phi }_s(Y)`$. Using 2.5 and 1.5, we get
$$\stackrel{~}{h}(Y,Y^{})=\stackrel{~}{h}(\mathrm{\Phi }_s(Y),Y^{})+_0^s\lambda _{\mathrm{\Phi }_\sigma (Y)}(V_G(\mathrm{\Phi }_\sigma (Y))G(\mathrm{\Phi }_\sigma (Y))d\sigma $$
$$h(q(s),t+s;q^{},t^{})+_0^sL(q(\sigma ),\dot{q}(\sigma ),t+\sigma )𝑑th(q,t;q^{},t^{})=\stackrel{~}{h}(Y,Y^{}).$$
As a consequence, all the inequalities are equalities. We obtain that the curve $`q(s)`$ is minimizing in the expression
$$h(q,t;q^{},t^{})=\mathrm{min}\left(h(q(s),t+s;q^{},t^{})+_0^sL(q(\sigma ),\dot{q}(\sigma ),t+\sigma )𝑑t\right).$$
Fathi has proved that $`P`$ is then a super-differential of the function $`h(.,t;q^{},t^{})`$ at $`q`$. The properties at $`(q^{},t^{})`$ are treated in a similar way.
### 3.4
Corollary If $`H`$ satisfies the hypotheses of 1.1, then $`m(H)\stackrel{~}{m}(H)`$.
### 3.5
Corollary If $`H`$ satisfies the hypotheses of 1.1, and if $`m(H)=0`$, then $`\stackrel{~}{m}(H)=0`$, and we have have $`\stackrel{~}{}𝒜_s(G)=\stackrel{~}{}𝒜(H)`$. In addition, we have
$$\stackrel{~}{h}(X_0,t_0,E_0;X_1,t_1,E_1)=h(\pi (P_0),t_0;\pi (P_1),t_1)$$
for each $`(P_0,t_0,E_0)`$ and $`(P_1,t_1,E_1)`$ in $`\stackrel{~}{}𝒜(H)`$.
Proof. Let $`(P,t,E)`$ be a point of $`T^{}N`$ and $`q=\pi (P)`$. If $`(P,t,E)\stackrel{~}{}𝒜_s(G)`$, then
$$\stackrel{~}{h}(P,t,E;P,t,E)=0,$$
so that $`h(q,t;q,t)0`$. Since, on the other hand, we have $`h(q,t;q,t)m(H)=0`$, we conclude that $`h(q,t;q,t)=0`$ hence $`(q,t)𝒜(H)`$. As a consequence, the function $`h(q,t;.,t)`$ is differentiable at $`q`$, see 1.6, and $`(_3h(q,t;q,t),tH(_3h(q,t;q,t),t))\stackrel{~}{}𝒜(H)`$. Since $`\stackrel{~}{h}(P,t,E;P,t,E)=h(q,t;q,t)`$, the point $`P`$ is a super-differential of $`h(q,t;.,t)`$ at $`q`$, and we must have $`P=_3h(q,t;q,t)`$. Moreover, we have $`G(P,t,E)=H(P,t)+E=0`$, hence $`(P,t,E)\stackrel{~}{}𝒜(H)`$.
Conversely, assume that $`(P,t,E)\stackrel{~}{}𝒜(H)`$. We then have $`E=H(P,t)`$. In addition, $`h(q,t;q,t)=0`$, the functions $`h(q,t;.,t)`$ and $`h(.,t;q,t)`$ are differentiable at $`q`$, and we have $`P=_3h(q,t;q,t)=_1h(q,t;q,t)`$. Now let $`XT_q^{}M`$ and $`X^{}T_q^{}M`$ be such that
$$\stackrel{~}{h}(X,t,H(X,t);X^{},t^{},H(X^{},t^{}))=h(q,t;q,t).$$
Then $`X`$ is a super-differential at $`q`$ of $`h(.,t;q,t)`$, and $`X^{}`$ is a super-differential at $`q`$ of $`h(q,t;.,t)`$. It follows that $`X=P=X^{}`$. Hence we have $`\stackrel{~}{h}(P,t,E;P,t,E)=h(q,t;q,t)=0`$. This proves that $`\stackrel{~}{m}(H)=0`$, and that $`(P,t,E)\stackrel{~}{}𝒜_s(G)`$.
Finally, let $`(P_0,t_0,E_0)T_{q_0}^{}M\times T^{}𝕋`$ and $`(P_1,t_1,E_1)T_{q_1}^{}M\times T^{}𝕋`$ be two points of $`\stackrel{~}{}𝒜(H)`$. We have $`E_0=H(P_0,t_0)`$ and $`E_1=H(P_1,t_1)`$. Furthermore, the function $`h(q_0,t_0;.,t_1)`$ is differentiable at $`q_1`$, with $`_3h(q_0,t_0;q_1,t_1)=P_1`$, and that the function $`h(.,t_0;q_1,t_1)`$ is differentiable at $`q_0`$, with $`_1h(q_0,t_0;q_1,t_1)=P_0`$. Since $`P_0`$ and $`P_1`$ are then the only super-differentials of $`h(.,t_0;q_1,t_1)`$ and $`h(q_0,t_0;.,t_1)`$, we conclude that $`\stackrel{~}{h}(P_0,t_0,E_0;P_1,t_1,E_1)=h(q_0,t_0;q_1,t_1)`$.
### 3.6
Corollary If $`H`$ satisfies the hypotheses of 1.1, and if $`m(H)=0`$, then $`\stackrel{~}{}_s(G)=\stackrel{~}{}(H)`$.
### 3.7
Corollary If $`H`$ satisfies the hypotheses of 1.1, and if $`m(H)=0`$, then $`\stackrel{~}{}𝒩_s(G)=\stackrel{~}{}𝒩(H)`$.
Proof. It is enough to prove that, if $`(P_0,t_0,E_0)`$ and $`(P_1,t_1,E_1)`$ belong to $`\stackrel{~}{}𝒜_s(G)`$, and $`q_0=\pi (P_0),q_1=\pi (P_1)`$, then
$$\stackrel{~}{}_s(P_0,t_0,E_0;P_1,t_1,E_1)=\stackrel{~}{}(q_0,t_0,q_1,t_1).$$
Let $`(P,t,E)`$ be a point of $`\stackrel{~}{}_s(P_0,t_0,E_0;P_1,t_1,E_1)`$. We then have $`G(P_0,t_0,E_0)=G(P,t,E)=0`$ hence $`E=H(P,t)`$. Furthermore, the inequalities
$$h(q_0,t_0;q_1,t_1)=\stackrel{~}{h}(P_0,t_0,H(P_0,t_0);P_1,t_1,H(P_1,t_1))$$
$$=\stackrel{~}{h}(P_0,t_0,H(P_0,t_1);P,t,E)+\stackrel{~}{h}(P,t,E;P_1,t_1,H(P_1,t_1))$$
$$h(q_0,t_0;q,t)+h(q,t;q_1,t_1)h(q_0,t_0;q_1,t_1)$$
are all equalities. As a consequence, the point $`(q,t)`$ belongs to the set $`(q_0,t_0;q_1,t_1)`$, and the differentials $`_3h(q_0,t_0;q,t)`$ and $`_1h(q,t;q_1,t_1)`$ exist, we have $`_3h(q_0,t_0;q,t)=_1h(q,t;q_1,t_1)`$, and the point
$$(X,t,e)=(_3h(q_0,t_0;q,t),t,H(_3h(q_0,t_0;q,t),t))$$
belongs to $`\stackrel{~}{}(q_0,t_0;q_1,t_1)`$, as follows from our definition of the Mañé set. Since
$$\stackrel{~}{h}(P_0,t_0,H(P_0,t_0);P,t,H(P,t))=h(q_0,t_0;q,t),$$
the point $`P`$ must be a super-differential of $`h(q_0,t_0;.,t)`$ at $`q`$, hence $`P=X`$. We have proved that $`(P,t,E)\stackrel{~}{}(q_0,t_0;q_1,t_1)`$.
Conversely, assume that $`(P,t,E)\stackrel{~}{}(q_0,t_0;q_1,t_1)`$, so that $`E=H(P,t)`$. Then
$$h(q_0,t_0;q,t)+h(q,t;q_1,t_1)=h(q_0,t_0;q_1,t_1)$$
and
$$P=_3h(q_0,t_0;q,t)=_1h(q,t;q_1,_1).$$
In addition, since $`(q_0,t_0)`$ and $`(q_1,t_1)`$ belong to $`𝒜(H)`$, the differential $`P_0=_1h(q_0,t_0;q,t)`$ exists for all $`q`$, and satisfies $`(P_0,t_0,H(P_0,t_0))\stackrel{~}{}𝒜(H))`$. Similarly, setting $`P_1=_3h(q,t;q_1,t_1)`$, we have $`(P_1,t_1,H(P_1,t_1))\stackrel{~}{}𝒜(H))`$. We conclude that
$$\stackrel{~}{h}(P_0,t_0,H(P_0,t_0);P,t,E)=h(q_0,t_0;q,t)$$
and
$$\stackrel{~}{h}(P,t,E;P_1,t_1,H(P_1,t_1))=h(q,t;q_1,t_1).$$
As a consequence, setting $`E_0=H(P_0,t_0)`$ and $`E_1=H(P_1,t_1)`$, we have
$$\stackrel{~}{h}(P_0,t_0,E_0;P,t,E)+\stackrel{~}{h}(P,t,E;P_1,t_1,E_1)=\stackrel{~}{h}(P_0,t_0,E_0;P_1,t_1,E_1).$$
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# Landau damping of Bogoliubov excitations in optical lattices at finite temperature
## I Introduction
Bose-Einstein condensates in a periodic optical lattice potential have attracted much interest. An optical lattice is an ideal system for studying a variety of solid state physics phenomena, such as Bloch oscillations Morsch ; Roati , Landau-Zener tunneling Anderson ; Jona-Lasinio , and Josephson oscillations Inguscio1 . An optical lattice is an ideal crystal without lattice imperfections or impurities, and parameters can be changed easily over a wide range. The experimental achievement of the superfluid-Mott insulator transition in an optical lattice Greiner also proved the usefulness of optical lattices to understand basic properties of strongly correlated systems. Moreover, Bose condensates in an optical lattice are predicted to show novel phenomena such as dynamical instability Wu2 ; Smerzi3 ; TaylorZaremba ; Fallani ; Modugno and swallow-tail energy loops in the band structure Wu1 ; Diakonov ; Machholm . Recent studies have concentrated on the case of pure Bose condensates at zero temperature. The effect of a thermal cloud on the dynamics of the condensate excitations in an optical lattice have not been studied very much. This is the subject of the present paper.
Recently, several experimental papers have reported results on the collective modes of Bose condensates in a one-dimensional optical lattice Inguscio3 ; Inguscio2 ; Esslinger2 ; Porto and its damping at finite temperature Inguscio2 . In these experiments, a shift of the oscillation frequency in the presence of the optical lattice and a sharp change of the damping rate with increasing depth of the optical lattice have been observed. The measured frequency shift Inguscio3 of the collective modes is in good agreement with the renormalized mass theory of Krämer et al. Kramer1 . The damping of condensate oscillations in a 1D optical lattice at $`T=0`$ has been measured Esslinger2 . The damping of the condensate excitations at finite $`T`$ has not been studied in any detail.
Since our major interest in this paper is the thermal cloud of excitations, it is important to make a clear distinction at the outset between (a) The Bloch-Bogoliubov excitations associated with linearized fluctuations of an equilibrium Bose-condensate and (b) The stationary states of the time-independent Gross-Pitaevskii equation for the Bose order parameter. In a continuum model, the latter states can be described by the eigenfunction (we use a 1D lattice for illustration)
$`\mathrm{\Phi }_k^0(x)=e^{\mathrm{i}kx}u_k(x),`$ (1)
where the condensate Bloch function satisfies the usual periodicity condition $`u_k(x)=u_k(x+ld)`$, where $`d`$ is the optical lattice spacing and $`l`$ is an integer. Physically, $`\mathrm{\Phi }_k^0(x)`$ corresponds to a solution of the GP equation with a superfluid flow in the periodic potential with the condensate quasi-momentum $`k`$. Recent theoretical literature (see, for example, Refs. Wu2 and TaylorZaremba ) has reported extensive studies of such condensate Bloch states, including their energy band structure and stability. The latter question can be studied by considering the dynamic fluctuations $`\delta \mathrm{\Phi }_k(x,q)`$ around the equilibrium state $`\mathrm{\Phi }_k^0(x)`$, with the generalized Bogoliubov excitation energy $`E_q(k)`$. These excitations are also described by a quasi-momentum $`q`$ in the first Brillouin zone (BZ) and will be referred to as the Bloch-Bogoliubov excitations of the optical lattice. The thermal cloud of non-condensate atoms which is present at finite temperature is described as a gas of these Bloch-Bogoliubov excitations. As emphasized in the literature Pitaevskii&Stringari , one must not confuse the energy bands of these Bloch-Bogoliubov excitations $`E_q(k)`$ with the condensate energy bands or Bloch eigenstates described by Eq. (1). That is, we must distinguish between the “condensate” energy band and the “excitation” energy band. In our tight-binding model, the analogue of Eq. (1) is
$`\mathrm{\Phi }_k^0(l)=e^{\mathrm{i}kld}\sqrt{n^\mathrm{c}(k)},`$ (2)
where $`d`$ is the optical lattice spacing and $`l`$ is an integer. While one could generalize our analysis, we only consider the Bose condensate in the $`k=0`$ Bloch state, in which case Eq. (2) reduces to
$`\mathrm{\Phi }_{k=0}^0(l)=\sqrt{n^{\mathrm{c0}}}.`$ (3)
Here $`n^{\mathrm{c0}}`$ denotes the number of condensate atoms trapped in each well of the optical lattice, labelled by $`l`$.
In this paper, we study the dynamics of Bose atoms in an optical lattice at finite temperatures. We discuss a Bose-Hubbard tight-binding model using the Gross-Pitaevskii approximation, but generalized to include the presence of non-condensate atoms. We calculate the temperature dependence of the condensate atom number $`n^{\mathrm{c0}}(T)`$ in each lattice well using the static Popov approximation. This is needed in our calculation of the Landau damping of condensate modes due to coupling to thermal excitations. We also extend our earlier results for a 1D optical lattice TsuchiyaGriffin to 2D and 3D optical lattices. As discussed in Ref. TsuchiyaGriffin , for damping processes to occur, the dispersion relation of the Bloch-Bogoliubov excitations $`E_q`$ must initially bend upward as the quasi-momentum $`q`$ increases. This is referred to as “anomalous dispersion” and is also the source of 3-phonon damping of long wavelength phonons in superfluid $`{}_{}{}^{4}\mathrm{He}`$ anomalous1 ; anomalous2 . This condition leads to a dramatic disappearance of all damping processes of phonon modes in a $`D`$-dimensional optical lattice when $`\alpha Un^{\mathrm{c0}}/J>6D`$, where $`U`$ is the on-site interaction and $`J`$ is the hopping matrix element.
This paper is organized as follows. In Sec. II we introduce the well-known Bose-Hubbard tight-binding model which describes Bose gases in an optical lattice. In Sec. III, we discuss the generalized discrete Gross-Pitaevskii equation and the characteristic changes in the Bloch-Bogoliubov excitation dispersion relation as a function of the dimensionless interaction parameter $`\alpha =Un^{\mathrm{c0}}/J`$. In Sec. IV, we introduce the static Popov approximation for optical lattices, and calculate the condensate fraction and $`\alpha `$ as a function of the optical lattice depth and temperature. In Sec. V, we calculate the Landau damping of Bogoliubov excitations in an optical lattice of dimension $`D`$. We also briefly remark on Beliaev damping into two excitations. We made some concluding remarks in Sec. VI
As mentioned, some of results of this paper in a 1D optical lattice were briefly described in Ref. TsuchiyaGriffin . In the present paper, we give a more detailed derivation and discussion, as well as new results for 2D and 3D optical lattices.
## II model
We consider bosonic atoms in an optical lattice potential
$`V_{\mathrm{op}}(𝐫)=sE_R{\displaystyle \underset{i=1}{\overset{D}{}}}\mathrm{sin}^2(kx_i),`$ (4)
where $`s`$ is the usual dimensionless parameter describing the optical lattice depth in units of the photon recoil energy $`E_R\mathrm{}^2k^2/2m`$. $`D`$ is the dimension of the optical lattice and $`d=\frac{\pi }{k}=\frac{\lambda }{2}`$ is the lattice period. We only consider simple cubic lattices considered in recent experiments Esslinger1 ; Greiner . We call attention to the recent technique Esslinger2 of producing a two-dimensional array of long, tightly confined condensate tubes by loading a Bose condensate into a deep 2D optical lattice potential, which prevents atoms from hopping between different tubes. With an additional 1D optical lattice potential along a tube, an ideal 1D system can be experimentally realized Esslinger2 . One can also have an ideal 2D system by loading a condensate into a deep 1D optical lattice and a shallow 2D optical lattice. We assume this experimental setup for the realization of 1D and 2D optical lattices in the present paper. We also assume that the laser intensity determining the depth of the optical lattice wells is large enough to make the atomic wave functions well localized on the individual sites (i.e., where we can use the tight-binding approximation). The energy gap between the first and the second excitation bands is large compared to the thermal energy ($`2k_\mathrm{B}T/E_Rs`$), and thus only the first band is thermally occupied.
Within a tight-binding approximation, the Hamiltonian is effectively described by the Bose-Hubbard model Fisher ; Jaksch ; Stoof ; ReyBurnett as
$`H=J{\displaystyle \underset{<j,l>}{}}(a_j^{}a_l+a_l^{}a_j)+{\displaystyle \frac{1}{2}}U{\displaystyle \underset{j}{}}a_j^{}a_j^{}a_ja_j,`$ (5)
where $`a_j`$ and $`a_j^{}`$ are destruction and creation operators of atoms on the $`j`$-th lattice site. $`j,l`$ represents nearest neighbor pairs of lattice sites. The first term describes the kinetic energy due to the hopping of atoms between sites. The hopping matrix element $`J`$ is given by
$`J={\displaystyle 𝑑𝐫w_j^{}(𝐫)\left(\frac{\mathrm{}^2^2}{2m}+V_{\mathrm{op}}(𝐫)\right)w_l(𝐫)},`$ (6)
where $`w_j(𝐫)`$ is a wave function localized on the $`j`$-th lattice site, and $`m`$ is the atomic mass. Expanding the optical lattice potential around the minima of the potential wells, the well trap frequency is $`\omega _ss^{1/2}\frac{\mathrm{}k^2}{m}`$. Approximating the localized function as the ground state wave function of a harmonic oscillator with frequency $`\omega _s`$ at the potential minima of $`j`$-th site
$`w_j(𝐫)=\left({\displaystyle \frac{m\omega _s}{\pi \mathrm{}}}\right)^{D/4}\mathrm{exp}\left({\displaystyle \frac{m\omega _s}{2\mathrm{}}}(𝐫𝐫_j)^2\right),`$ (7)
one obtains
$`{\displaystyle \frac{J}{E_R}}\left[({\displaystyle \frac{\pi ^2s}{4}}{\displaystyle \frac{s^{1/2}}{2}}){\displaystyle \frac{1}{2}}s(1+\mathrm{exp}(s^{1/2}))\right]e^{\pi ^2s^{1/2}/4}.`$ (8)
Here $`s^{1/2}\frac{\sqrt{2m(sE_R)}}{\mathrm{}}d`$ can be interpreted as a WKB factor for tunneling in an optical lattice potential which has height $`sE_R`$ and well width $`d`$.
The second term in Eq. (5) describes the interaction between atoms when they are at the same site. We assume that atoms can move along $`z`$ direction in 1D case and in $`xy`$ plane in 2D case. The on-site interaction $`U`$ depends on the dimensionality of the optical lattice. $`U`$ is given by Jaksch
$`U=g{\displaystyle d𝐫_{}\left|\varphi _{}(𝐫_{})\right|^4dz|w_j(z)|^4},`$ $`(1\mathrm{D})`$ (9)
$`=g{\displaystyle dz\left|\varphi _{}(z)\right|^4dxdy|w_j(x,y)|^4},`$ $`(2\mathrm{D})`$ (10)
$`=g{\displaystyle d𝐫|w_j(𝐫)|^4},`$ $`(3\mathrm{D})`$ (11)
where $`g=4\pi \mathrm{}^2a/m`$ and $`a`$ is the $`s`$-wave scattering length. Here $`\varphi _{}(𝐫_{})=\left(\frac{m\omega _{}}{\pi \mathrm{}}\right)^{1/2}\mathrm{exp}\left(\frac{m\omega _{}}{2\mathrm{}}𝐫_{}^2\right)`$ and $`\varphi _{}(z)=\left(\frac{m\omega _{}}{\pi \mathrm{}}\right)^{1/4}\mathrm{exp}\left(\frac{m\omega _{}}{2\mathrm{}}z^2\right)`$ are the ground state wave functions in optical lattice well traps for confining atoms in 1 and 2 dimensions. Approximating the localized function $`w_j`$ as a simple Gaussian, one obtains Jaksch
$`{\displaystyle \frac{U}{E_R}}{\displaystyle \frac{g}{(2\pi )^{3/2}a_{}^2a_s}}={\displaystyle \frac{2^{3/2}ad}{\pi ^{3/2}a_{}^2}}s^{1/4},`$ $`(1\mathrm{D})`$ (12)
$`{\displaystyle \frac{g}{(2\pi )^{3/2}a_{}a_s^2}}={\displaystyle \frac{2^{3/2}a}{\pi ^{1/2}a_{}}}s^{1/2},`$ $`(2\mathrm{D})`$ (13)
$`{\displaystyle \frac{g}{(2\pi )^{3/2}a_s^3}}={\displaystyle \frac{2^{3/2}\pi ^{1/2}a}{d}}s^{3/4},`$ $`(3\mathrm{D})`$ (14)
where $`a_s=\sqrt{\frac{\mathrm{}}{m\omega _s}}`$, $`a_{}=\sqrt{\frac{\mathrm{}}{m\omega _{}}}`$, and $`a_{}=\sqrt{\frac{\mathrm{}}{m\omega _{}}}`$.
## III the generalized Gross-Pitaevskii equation
In order to investigate the collective excitations of a Bose condensate in an optical lattice, we restrict ourselves to the superfluid phase of the system described by Eq. (5). We start from the Heisenberg equation of motion for $`a_j(t)`$ (we set $`\mathrm{}=1`$ from now on),
$`\mathrm{i}{\displaystyle \frac{\mathrm{d}}{\mathrm{d}t}}a_j(t)`$ $`=`$ $`[a_j(t),H]`$ (15)
$`=`$ $`J{\displaystyle \underset{l}{}}a_l(t)+Ua_j^{}(t)a_j(t)a_j(t),`$ (16)
where $`l`$ means that $`l`$ runs over the nearest neighbor sites of site $`j`$. In the presence of Bose condensation, we write $`a_j(t)=\mathrm{\Phi }_j(t)+\stackrel{~}{\psi }_j(t)`$, where $`\mathrm{\Phi }_j(t)=a_j(t)`$ is the condensate wave function at site $`j`$ and $`\stackrel{~}{\psi }_j(t)`$ is the non-condensate field operator (with $`\stackrel{~}{\psi }_j(t)=0`$). This non-condensate operator $`\stackrel{~}{\psi }_j`$ satisfies the usual Bose commutation relations,
$`[\stackrel{~}{\psi }_j,\stackrel{~}{\psi }_l^{}]=\delta _{jl},[\stackrel{~}{\psi }_j,\stackrel{~}{\psi }_l]=[\stackrel{~}{\psi }_j^{},\stackrel{~}{\psi }_l^{}]=0.`$ (17)
The equation for the condensate wave function is obtained by taking an average of the Heisenberg equation. If, as usual, we neglect the anomalous correlations $`m_j(t)=\stackrel{~}{\psi }_j\stackrel{~}{\psi }_j`$ and $`\stackrel{~}{\psi }_j^{}\stackrel{~}{\psi }_j\stackrel{~}{\psi }_j`$, we are left with a “generalized” form of the discrete Gross-Pitaevskii equation Hutchinson which includes the mean field of non-condensate atoms
$`\mathrm{i}{\displaystyle \frac{\mathrm{d}}{\mathrm{d}t}}\mathrm{\Phi }_j(t)=J{\displaystyle \underset{l}{}}\mathrm{\Phi }_l(t)+U(n_j^\mathrm{c}(t)+2\stackrel{~}{n}_j(t))\mathrm{\Phi }_j(t).`$ (18)
Here, $`n_j^\mathrm{c}(t)`$ is the number of condensate atoms on the $`j`$-th lattice site and $`\stackrel{~}{n}_j(t)=\stackrel{~}{\psi }_j^{}\stackrel{~}{\psi }_j`$ is the number of the non-condensate atoms on $`j`$-th lattice site. The time-dependent Hartree-Fock mean field $`2U\stackrel{~}{n}_j(t)`$ in Eq. (18) arises from the thermal gas of non-condensate atoms. Equation (18) reduces to the usual discrete Gross-Pitaevskii equation Inguscio1 ; Smerzi3 ; ReyBurnett when all the atoms are assumed to be in the condensate (i.e., if we set $`\stackrel{~}{n}_j=0`$).
Introducing the usual phase and amplitude variables for the condensate, $`\mathrm{\Phi }_j(t)=\sqrt{n_j^\mathrm{c}(t)}e^{\mathrm{i}\theta _j(t)}`$, Eq. (18) is equivalent to
$`{\displaystyle \frac{\mathrm{d}n_j^\mathrm{c}(t)}{\mathrm{d}t}}`$ $`=`$ $`2J{\displaystyle \underset{l}{}}\sqrt{n_j^\mathrm{c}(t)n_l^\mathrm{c}(t)}\mathrm{sin}(\theta _l(t)\theta _j(t)),`$ (19)
$`{\displaystyle \frac{\mathrm{d}\theta _j(t)}{\mathrm{d}t}}`$ $`=`$ $`J{\displaystyle \underset{l}{}}\sqrt{{\displaystyle \frac{n_l^\mathrm{c}(t)}{n_j^\mathrm{c}(t)}}}\mathrm{cos}(\theta _l(t)\theta _j(t))U(n_j^\mathrm{c}(t)+2\stackrel{~}{n}_j(t))`$ (20)
$``$ $`\epsilon _j^\mathrm{c}(t).`$
Equation (19) is the continuity equation for the condensate and Eq. (20) often called the Josephson equation. The energy of a condensate atom $`\epsilon _j^\mathrm{c}(t)`$ reduces to the equilibrium chemical potential $`\mu _{\mathrm{c0}}`$ of the condensate in static thermal equilibrium ($`\theta _j`$, $`n_j^\mathrm{c}`$ and $`\stackrel{~}{n}_j`$ are independent of site index $`j`$)
$`\mu _{\mathrm{c0}}=zJ+U(n^{\mathrm{c0}}(T)+2\stackrel{~}{n}^0(T)).`$ (21)
Here $`z`$ is the number of the nearest neighbor sites, $`n^{\mathrm{c0}}`$ is the number of the condensate atoms per site in equilibrium, and $`\stackrel{~}{n}^0`$ is the number of the non-condensate atoms per site in equilibrium. The solution of Eq. (20) in static equilibrium is $`\theta _0(t)=\mu _{\mathrm{c0}}t`$. From Eq. (19), one finds that the Josephson current between the $`j`$-th and $`l`$-th lattice sites is $`J_j(t)=2J\sqrt{n_j^\mathrm{c}(t)n_l^\mathrm{c}(t)}\mathrm{sin}(\theta _l(t)\theta _j(t))`$.
The well-known Bogoliubov excitation spectrum for a uniform optical lattice is easily obtained from the above equations Stoof ; ReyBurnett ; Javanainen . Considering small fluctuations from equilibrium, $`n_j^\mathrm{c}(t)=n^{\mathrm{c0}}+\delta n_j^\mathrm{c}(t)`$, $`\theta _j(t)=\theta _0(t)+\delta \theta _j(t)`$ and $`\stackrel{~}{n}_j(t)=\stackrel{~}{n}^0+\delta \stackrel{~}{n}_j(t)`$, Eqs. (19) and (20) reduce to
$`{\displaystyle \frac{\mathrm{d}\delta n_j^\mathrm{c}(t)}{\mathrm{d}t}}`$ $`=`$ $`2Jn^{\mathrm{c0}}{\displaystyle \underset{l}{}}\left(\delta \theta _l(t)\delta \theta _j(t)\right),`$ (22)
$`{\displaystyle \frac{\mathrm{d}\delta \theta _j(t)}{\mathrm{d}t}}`$ $`=`$ $`{\displaystyle \frac{J}{2n^{\mathrm{c0}}}}{\displaystyle \underset{l}{}}\left(\delta n_l^\mathrm{c}(t)\delta n_j^\mathrm{c}(t)\right)U\left(\delta n_j^\mathrm{c}(t)+2\delta \stackrel{~}{n}_j(t)\right).`$ (23)
Ignoring the non-condensate atom term ($`\delta \stackrel{~}{n}_j=0`$), the solution of these coupled equations are the normal modes
$`\delta \theta _j(t)=\delta \theta (𝐪)e^{\mathrm{i}[𝐪𝐫_jE_𝐪t]},\delta n_j^\mathrm{c}(t)=\delta n^\mathrm{c}(𝐪)e^{\mathrm{i}[𝐪𝐫_jE_𝐪t]}.`$ (24)
The Bloch-Bogoliubov excitation energy in an optical lattice is given by
$`E_𝐪`$ $`=`$ $`\sqrt{ϵ_𝐪^0(ϵ_𝐪^0+2Un^{\mathrm{c0}})},`$ (25)
where
$`ϵ_𝐪^0`$ $``$ $`4J{\displaystyle \underset{i=1}{\overset{D}{}}}\mathrm{sin}^2{\displaystyle \frac{q_id}{2}}.`$ (26)
For small $`q`$, this spectrum is phonon-like $`E_𝐪cq`$, with the phonon velocity
$`c=\sqrt{2Jd^2Un^{\mathrm{c0}}}=\sqrt{{\displaystyle \frac{Un^{\mathrm{c0}}}{m^{}}}},`$ (27)
where $`m^{}=\frac{1}{2Jd^2}`$ is an effective mass of atoms in the first excitation band of the optical lattice Kramer1 ; Kramer2 . This $`T=0`$ excitation spectrum in 1D case is shown in Fig. 1, for two values of the dimensionless interaction parameter $`\alpha Un^{\mathrm{c0}}/J`$.
We call attention to an important feature of the dispersion relation $`E_q`$ in Fig. 1, considered as a function of the ratio $`\alpha `$. For $`\alpha 6`$, the excitation energy $`E_q`$ bends up before bending over, as $`q`$ approaches the BZ boundary. This behavior is analogous to the so called “anomalous dispersion” of the phonon spectrum in superfluid $`{}_{}{}^{4}\mathrm{He}`$ anomalous1 ; anomalous2 ; Pitaevskii&Levinson . For $`\alpha >6`$, in contrast, the spectrum simply bends over as one leaves the low $`q`$ (phonon) region. This feature will play a crucial role when we consider damping processes in optical lattices. The excitation spectrum in 2D and 3D optical lattices also exhibit this kind of spectrum. However, the critical value of $`\alpha `$ then depends on the direction of $`𝐪`$, since simple cubic optical lattices in 2D and 3D do not have rotational symmetry. The crucial effect of this anomalous dispersion on damping processes will be discussed in detail in section V.
## IV static popov approximation
In this section, we calculate the condensate fraction $`n^{\mathrm{c0}}/n`$ and the parameter $`\alpha Un^{\mathrm{c0}}/J`$ as a function of the optical lattice depth $`s`$ and the temperature. This is needed as input to the calculations in Sec. V. We base our discussion on the static Popov approximation Hutchinson ; ReyBurnett ; WG .
Substituting $`a_j=\mathrm{\Phi }_j+\stackrel{~}{\psi }_j`$ to $`KH\mu N`$, where $`N=_ja_j^{}a_j`$ is the total number of atoms, one obtains,
$`K`$ $`=`$ $`K_0+K_1+K_2+K_3+K_4,`$ (28)
$`K_0`$ $`=`$ $`J{\displaystyle \underset{j,l}{}}\left(\mathrm{\Phi }_j^{}\mathrm{\Phi }_l+\mathrm{\Phi }_l^{}\mathrm{\Phi }_j\right)\mu {\displaystyle \underset{j}{}}|\mathrm{\Phi }_j|^2+{\displaystyle \frac{U}{2}}{\displaystyle \underset{j}{}}|\mathrm{\Phi }_j|^4,`$ (29)
$`K_1`$ $`=`$ $`{\displaystyle \underset{j}{}}\left(J{\displaystyle \underset{l}{}}\mathrm{\Phi }_l\mu \mathrm{\Phi }_j+U|\mathrm{\Phi }_j|^2\mathrm{\Phi }_j\right)\stackrel{~}{\psi }_j^{}`$ (30)
$`+\left(J{\displaystyle \underset{l}{}}\mathrm{\Phi }_l^{}\mu \mathrm{\Phi }_j^{}+U|\mathrm{\Phi }_j|^2\mathrm{\Phi }_j^{}\right)\stackrel{~}{\psi }_j`$
$`K_2`$ $`=`$ $`J{\displaystyle \underset{j,l}{}}\left(\stackrel{~}{\psi }_j^{}\stackrel{~}{\psi }_l+\stackrel{~}{\psi }_l^{}\stackrel{~}{\psi }_j\right)\mu {\displaystyle \underset{j}{}}\stackrel{~}{\psi }_j^{}\stackrel{~}{\psi }_j`$ (31)
$`+{\displaystyle \frac{U}{2}}{\displaystyle \underset{j}{}}\left(\mathrm{\Phi }_j^2(\stackrel{~}{\psi }_j^{})^2+4|\mathrm{\Phi }_j|^2\stackrel{~}{\psi }_j^{}\stackrel{~}{\psi }_j+(\mathrm{\Phi }_j^{})^2\stackrel{~}{\psi }_j^2\right),`$
$`K_3`$ $`=`$ $`U{\displaystyle \underset{j}{}}\left(\mathrm{\Phi }_j^{}\stackrel{~}{\psi }_j^{}\stackrel{~}{\psi }_j^2+\mathrm{\Phi }_j(\stackrel{~}{\psi }_j^{})^2\stackrel{~}{\psi }_j\right),`$ (32)
$`K_4`$ $`=`$ $`{\displaystyle \frac{U}{2}}{\displaystyle \underset{j}{}}(\stackrel{~}{\psi }_j^{})^2\stackrel{~}{\psi }_j^2.`$ (33)
Here, we use the Hartree-Fock-Bogoliubov-Popov approximation that takes into account the third and fourth order terms of $`\stackrel{~}{\psi }_j`$, $`\stackrel{~}{\psi }_j^{}`$ within a mean-field approximation, but neglects terms involving the anomalous averages $`\stackrel{~}{\psi }_j\stackrel{~}{\psi }_j`$ and $`\stackrel{~}{\psi }_j^{}\stackrel{~}{\psi }_j\stackrel{~}{\psi }_j`$, namely we use
$`K_32U{\displaystyle \underset{j}{}}\stackrel{~}{n}_j\left(\mathrm{\Phi }_j^{}\stackrel{~}{\psi }_j+\mathrm{\Phi }_j\stackrel{~}{\psi }_j^{}\right),K_42U{\displaystyle \underset{j}{}}\stackrel{~}{n}_j\stackrel{~}{\psi }_j^{}\stackrel{~}{\psi }_j.`$ (34)
Within this approximation, we have
$`K`$ $`=`$ $`K_0+K_1^{}+K_2^{}`$ (35)
$`K_1^{}`$ $`=`$ $`{\displaystyle \underset{j}{}}\left(J{\displaystyle \underset{l}{}}\mathrm{\Phi }_l\mu \mathrm{\Phi }_j+U(n_j^\mathrm{c}+2\stackrel{~}{n}_j)\mathrm{\Phi }_j\right)\stackrel{~}{\psi }_j^{}`$ (36)
$`+\left(J{\displaystyle \underset{l}{}}\mathrm{\Phi }_l^{}\mu \mathrm{\Phi }_j^{}+U(n_j^\mathrm{c}+2\stackrel{~}{n}_j)\mathrm{\Phi }_j^{}\right)\stackrel{~}{\psi }_j`$
$`K_2^{}`$ $`=`$ $`J{\displaystyle \underset{j,l}{}}\left(\stackrel{~}{\psi }_j^{}\stackrel{~}{\psi }_l+\stackrel{~}{\psi }_l^{}\stackrel{~}{\psi }_j\right)\mu {\displaystyle \underset{j}{}}\stackrel{~}{\psi }_j^{}\stackrel{~}{\psi }_j`$ (37)
$`+{\displaystyle \frac{U}{2}}{\displaystyle \underset{j}{}}\left(\mathrm{\Phi }_j^2(\stackrel{~}{\psi }_j^{})^2+4n_j\stackrel{~}{\psi }_j^{}\stackrel{~}{\psi }_j+(\mathrm{\Phi }_j^{})^2\stackrel{~}{\psi }_j^2\right),`$
where $`n_jn_j^\mathrm{c}+\stackrel{~}{n}_j=|\mathrm{\Phi }_j|^2+\stackrel{~}{\psi }_j^{}\stackrel{~}{\psi }_j`$. The coefficients of the linear terms in $`\stackrel{~}{\psi }`$, $`\stackrel{~}{\psi }^{}`$ in Eq. (36) identically vanish using the fact that $`\mathrm{\Phi }_j`$ satisfies the generalized Gross-Pitaevskii equation Eq. (18). As a result, $`K_1^{}=0`$.
As discussed in Sec. I, we consider a Bose condensate in static thermal equilibrium in the $`k=0`$ Bloch state, and hence $`\mathrm{\Phi }_j=\sqrt{n^{\mathrm{c0}}}`$. Using the value of the condensate chemical potential in thermal equilibrium given in Eq. (21), the remaining term $`K_2^{}`$ in Eq. (36) reduces to
$`K_2^{}`$ $`=`$ $`J{\displaystyle \underset{j,l}{}}\left(\stackrel{~}{\psi }_j^{}\stackrel{~}{\psi }_l+\stackrel{~}{\psi }_l^{}\stackrel{~}{\psi }_j\right)+(zJ+Un^{\mathrm{c0}}){\displaystyle \underset{j}{}}\stackrel{~}{\psi }_j^{}\stackrel{~}{\psi }_j`$ (38)
$`+{\displaystyle \frac{U}{2}}{\displaystyle \underset{j}{}}\left(n^{\mathrm{c0}}(\stackrel{~}{\psi }_j^{})^2+n^{\mathrm{c0}}\stackrel{~}{\psi }_j^2\right).`$
We next introduce Fourier components,
$`\stackrel{~}{\psi }_j={\displaystyle \frac{1}{\sqrt{I^D}}}{\displaystyle \underset{𝐪0}{}}a_𝐪e^{\mathrm{i}𝐪𝐫_j},`$ (39)
where the number of lattice sites in one direction is denoted by $`I`$ and hence the total number of sites is $`I^D`$. The Momentum sum is over the first Brilloiuin zone of the optical lattice. Substituting Eq. (39) into Eq.(38), one finds
$`K_2^{}={\displaystyle \underset{𝐪0}{}}(ϵ_𝐪^0+Un^{\mathrm{c0}})a_𝐪^{}a_𝐪+{\displaystyle \frac{Un^{\mathrm{c0}}}{2}}{\displaystyle \underset{𝐪0}{}}(a_𝐪^{}a_𝐪^{}+a_𝐪a_𝐪).`$ (40)
We can diagonalize $`K=K_0+K_2^{}`$ by the Bogoliubov transformation Stoof ; ReyBurnett ,
$`\alpha _𝐪`$ $`=`$ $`u_𝐪a_𝐪+v_𝐪a_𝐪^{},`$
$`\alpha _𝐪^{}`$ $`=`$ $`v_𝐪^{}a_𝐪+u_𝐪^{}a_𝐪^{}.`$ (41)
If we assume that $`\alpha _𝐪`$ and $`\alpha _𝐪^{}`$ obey the usual Bose commutation relations, we obtain the following conditions for $`u_𝐪`$ and $`v_𝐪,`$
$`|u_𝐪|^2|v_𝐪|^2=1.`$ (42)
From the condition for the diagonalization of $`K`$, we find that $`u_𝐪`$ and $`v_𝐪`$ have to satisfy the following equation,
$`Un^{\mathrm{c0}}(u_𝐪^2+v_𝐪^2)2(ϵ_𝐪^0+Un^{\mathrm{c0}})u_𝐪v_𝐪=0.`$ (43)
Solving both Eq. (42) and Eq. (43), one can easily derive the parameters for the Bogoliubov transformation (we take $`u_𝐪`$ and $`v_𝐪`$ real),
$`u_𝐪^2={\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\stackrel{~}{E}_𝐪}{E_𝐪}}+1\right),v_𝐪^2={\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\stackrel{~}{E}_𝐪}{E_𝐪}}1\right).`$ (44)
We have here introduced the Hartree-Fock (HF) excitation spectrum
$`\stackrel{~}{E}_𝐪`$ $``$ $`[ϵ_0(q)+2U(n^{\mathrm{c0}}+\stackrel{~}{n}^0)]\mu _{\mathrm{c0}}`$ (45)
$`=`$ $`4J\mathrm{sin}^2({\displaystyle \frac{qd}{2}})+Un^{\mathrm{c0}}.`$
Putting all these results together, our total Hamiltonian has been reduced to
$`K`$ $`=`$ $`K_0+{\displaystyle \underset{𝐪0}{}}\left(\stackrel{~}{E}_𝐪v_𝐪^2Un^{\mathrm{c0}}u_𝐪v_𝐪\right)+{\displaystyle \underset{𝐪0}{}}E_𝐪\alpha _𝐪^{}\alpha _𝐪`$ (46)
$`=`$ $`{\displaystyle \frac{UN_{\mathrm{c0}}}{2}}(n^{\mathrm{c0}}+4\stackrel{~}{n}^0)+{\displaystyle \frac{1}{2}}{\displaystyle \underset{𝐪0}{}}\left(E_𝐪\stackrel{~}{E}_𝐪\right)+{\displaystyle \underset{𝐪0}{}}E_𝐪\alpha _𝐪^{}\alpha _𝐪.`$
Here $`N_{\mathrm{c0}}=n^{\mathrm{c0}}I^D`$ is the total number of condensate atoms. The Bogoliubov-Popov excitation spectrum $`E_𝐪`$ appearing in Eq. (46) is identical to Eq. (25), except that now $`n^{\mathrm{c0}}(T)`$ is the temperature-dependent number of condensate atoms in any given lattice well. The Bogoliubov-Popov excitation spectrum can be obtained by ignoring the non-condensate fluctuation ($`\delta \stackrel{~}{n}_j=0`$) but using the temperature-dependent number of condensate atoms $`n^{\mathrm{c0}}(T)`$ as given by Eq. (22) and Eq. (23). We thus assume that the thermal cloud is always in static equilibrium when dealing with the time-dependent density fluctuations $`\delta n_j^\mathrm{c}(t)`$ of the condensate. The approximation of a static thermal cloud has been successfully used in other discussions of condensate collective modes WG ; Hutchinson , which we refer the reader for further discussion.
It is easy to verify that $`E_𝐪`$ in Eq. (25) reduces to the HF energy $`\stackrel{~}{E}_𝐪`$ in the limit $`JUn^{\mathrm{c0}}`$, i.e., when $`\alpha 1`$. This HF limit corresponds to setting the Bogoliubov amplitudes $`u_𝐪^2=1`$, $`v_𝐪^2=0`$ in Eq. (44). In dealing with Bose gases trapped in harmonic potentials, one can always use this HF approximation ZNG for the excitations describing the thermal cloud as long as the kinetic energy of the atoms ($`k_\mathrm{B}T`$) is much larger than the interaction energy ($`Un^{\mathrm{c0}}`$). In contrast, apart from the limiting case of $`\alpha 1`$, we must always use the full Bogoliubov spectrum $`E_𝐪`$ to describe the thermal cloud composed of excitations in the first band of an optical lattice.
Expressing $`\stackrel{~}{n}^0(T)`$ in terms of these Bogoliubov-Popov excitations, we have ReyBurnett ; Stoof ; Hutchinson
$`n`$ $`=`$ $`n^{\mathrm{c0}}+{\displaystyle \frac{1}{I^D}}{\displaystyle \underset{𝐪0}{}}a_𝐪^{}a_𝐪`$ (47)
$`=`$ $`n^{\mathrm{c0}}+{\displaystyle \frac{1}{I^D}}{\displaystyle \underset{𝐪0}{}}\left[\left(u_𝐪^2+v_𝐪^2\right)f^0(E_𝐪)+v_𝐪^2\right],`$
where $`f^0(E_q)=[\mathrm{exp}(\beta E_q)1]^1`$ is the usual Bose distribution function. The number of condensate atoms $`n^{\mathrm{c0}}`$ at a site is found by solving Eq. (47) self-consistently for a fixed value of the total site density $`n`$. The condensate fraction $`n^{\mathrm{c0}}/n`$ in a $`D`$-dimensional optical lattice is shown in Fig. 2. We take $`n=2`$ and use the parameters from Ref. Esslinger2 The spurious finite jump in the condensate atom number $`n^{\mathrm{c0}}`$ at the transition temperature $`T_\mathrm{c}`$ is an inherent problem of the Bogoliubov theory in a uniform gas (for further discussion see, for example, Ref. Shi ).
Strictly speaking, there are no solutions of Eq. (47) at finite temperature for an infinite optical lattice in 1D and 2D because of the divergent contribution from excitations with small momentum, in accordance with the well-known Mermin-Wagner-Hohenberg theorem Mermin\_Wagner . However, for the finite systems discussed in this paper, this divergence is not present. That is, Eq. (47) has a solution describing a finite value of the condensate $`n^{\mathrm{c0}}(T)`$ below $`T_\mathrm{c}`$.
In Fig. 3, we plot the parameter $`\alpha Un^{\mathrm{c0}}/J`$ as a function of the temperature for several values of the optical depth $`s`$. The dimensionless interaction parameter $`\alpha `$ and the results in Fig. 3 will play an important role in the discussion in the rest of this paper. Since we limit our discussion to the first energy band of the optical lattice, our results only apply when $`s2k_\mathrm{B}T/E_\mathrm{R}`$. Higher excitation bands would be thermally populated if we consider lower values of $`s`$ and would have to be considered.
## V damping of Bogoliubov excitations
In this section, we discuss the damping of condensate modes in optical lattices. In Sec. V.1, we derive a formal expression for the damping of Bogoliubov excitations using the thermal Green’s function formalism. In Sec. V.2, we use this expression to calculate the Landau damping of Bogoliubov excitations in 1D lattice. Section V.3 extends this analysis to 2D and 3D optical lattices.
### V.1 Green’s function technique
In this section, we take $`K_\mathrm{B}K_0+K_1+K_2`$ in Eq. (33) as the bare zeroth-order Hamiltonian and treat $`K_3`$ as a perturbation (We neglect the fourth order term $`K_4`$). We use the simplest Bogoliubov theory to diagonalize $`K_\mathrm{B}`$ Stoof ; ReyBurnett . We neglect the fourth order term in the fluctuation $`K_4`$. The chemical potential for $`K_B`$ is given by $`\mu _{\mathrm{B0}}=zJ+Un^{\mathrm{c0}}`$ (from the condition $`K_1=0`$). $`K_2`$ can be diagonalized by the Bogoliubov transformation Eq. (41) and one obtains Eq. (46).
Using Eq. (39) and Eq. (41), $`K_3`$ can be re-written in terms of the Bogoliubov quasiparticles
$`K_3`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{𝐩_1,𝐩_2,𝐩_30}{}}M_{𝐩_1,𝐩_2;𝐩_3}\left(\alpha _{𝐩_3}^{}\alpha _{𝐩_1}\alpha _{𝐩_2}+\alpha _{𝐩_3}\alpha _{𝐩_1}^{}\alpha _{𝐩_2}^{}\right)`$ (48)
$`+{\displaystyle \underset{𝐩_1,𝐩_2,𝐩_30}{}}L_{𝐩_1,𝐩_2;𝐩_3}\left(\alpha _{𝐩_3}^{}\alpha _{𝐩_1}^{}\alpha _{𝐩_2}^{}+\alpha _{𝐩_3}\alpha _{𝐩_1}\alpha _{𝐩_2}\right),`$
where the matrix elements are
$`M_{𝐩_1,𝐩_2;𝐩_3}`$ $`=`$ $`2U\sqrt{{\displaystyle \frac{n^{\mathrm{c0}}}{I^D}}}{\displaystyle \underset{𝐆}{}}[(u_{𝐩_1}u_{𝐩_3}+v_{𝐩_1}v_{𝐩_3}v_{𝐩_1}u_{𝐩_3})u_{𝐩_2}`$ (49)
$`(u_{𝐩_1}u_{𝐩_3}+v_{𝐩_1}v_{𝐩_3}u_{𝐩_1}v_{𝐩_3})v_{𝐩_2}]\delta _{𝐩_1+𝐩_2,𝐩_3+𝐆},`$
$`L_{𝐩_1,𝐩_2;𝐩_3}`$ $`=`$ $`U\sqrt{{\displaystyle \frac{n^{\mathrm{c0}}}{I^D}}}{\displaystyle \underset{𝐆}{}}\left(v_{𝐩_1}v_{𝐩_2}u_{𝐩_3}u_{𝐩_1}u_{𝐩_2}v_{𝐩_3}\right)\delta _{𝐩_1+𝐩_2,𝐩_3+𝐆},`$ (50)
and $`𝐆`$ is a reciprocal lattice vector. The Kronecker delta $`\delta _{𝐩_1+𝐩_2,𝐩_3+𝐆}`$ expresses the conservation of quasi-momentum of the three excitation scattering processes (Umklapp processes are associated with $`𝐆0`$).
The thermal Green’s function for Bogoliubov quasiparticles is defined by (see e.g. Mahan )
$`G_𝐪(\tau _1\tau _2)T_\tau \left(\alpha _𝐪(\tau _1)\alpha _𝐪^{}(\tau _2)\right),`$ (51)
where
$`\alpha _𝐪(\tau )e^{K\tau }\alpha _𝐪e^{K\tau },\alpha _𝐪^{}(\tau )e^{K\tau }\alpha _𝐪^{}e^{K\tau }.`$ (52)
The angular brackets $`\mathrm{}`$ indicates taking an average for thermal equilibrium and $`T_\tau `$ gives the usual time ordered product of operators. Fourier transformation of this imaginary time Green’s function is given by
$`G_𝐪(\tau )={\displaystyle \frac{1}{\beta }}{\displaystyle \underset{\omega _n}{}}G_𝐪(\mathrm{i}\omega _n)e^{\mathrm{i}\omega _n\tau },`$ (53)
where $`\omega _n`$ is the boson Matsubara frequency $`\omega _n\frac{2n\pi }{\beta }`$ ($`n`$ is an integer). The zero-th order Green’s function is
$`G_𝐪^0(\mathrm{i}\omega _n)={\displaystyle \frac{1}{\mathrm{i}\omega _nE_𝐪}},`$ (54)
where $`E_𝐪`$ is the Bloch-Bogoliubov excitation energy given in Eq. (25).
The second-order self-energy terms are given by the usual diagrams shown in Fig. 4.
Fig. 4 (a) describes a condensate excitation of (quasi) momentum $`𝐪`$ being absorbed by an excitation $`𝐩_2`$ of the optical lattice thermal gas, lending to a thermal excitation with momentum $`𝐩_1`$. This describes Landau damping processes Szepfalusy ; PitaevskiiStringari ; Fedichev . Fig. 4 (b) describes a condensate excitation of momentum $`𝐪`$ decay into two excitations with momentums $`𝐩_1`$ and $`𝐩_2`$. This describes Beliaev damping Beliaev . These diagrams give the second-order self-energy of the Bogoliubov excitation
$`\mathrm{\Sigma }_𝐪(\mathrm{i}\omega _n)`$ $`=`$ $`\mathrm{\Sigma }_𝐪^\mathrm{L}(\mathrm{i}\omega _n)+\mathrm{\Sigma }_𝐪^\mathrm{B}(\mathrm{i}\omega _n),`$
$`\mathrm{\Sigma }_𝐪^\mathrm{L}(\mathrm{i}\omega _n)`$ $`=`$ $`{\displaystyle \frac{1}{\beta }}{\displaystyle \underset{𝐩_1,𝐩_20}{}}{\displaystyle \underset{\omega _l}{}}|M_{𝐪,𝐩_2;𝐩_1}|^2G_{𝐩_1}^0(\mathrm{i}\omega _l)G_{𝐩_2}^0(\mathrm{i}\omega _l\mathrm{i}\omega _n),`$
$`\mathrm{\Sigma }_𝐪^\mathrm{B}(\mathrm{i}\omega _n)`$ $`=`$ $`{\displaystyle \frac{1}{2\beta }}{\displaystyle \underset{𝐩_1,𝐩_20}{}}{\displaystyle \underset{\omega _l}{}}|M_{𝐩_1,𝐩_2;𝐪}|^2G_{𝐩_1}^0(\mathrm{i}\omega _l)G_{𝐩_2}^0(\mathrm{i}\omega _n\mathrm{i}\omega _l).`$ (55)
The summation over Matsubara frequencies is calculated by the well-known formula (see Ref. Mahan )
$`{\displaystyle \underset{\omega _n}{}}g(\omega _n)={\displaystyle \frac{\beta }{2\pi \mathrm{i}}}{\displaystyle _C}dzg(z)f^0(z),`$ (56)
where $`f^0(z)`$ is the Bose distribution function and the contour encircles the poles of $`g(z)`$. Using Eq. (56), one obtains
$`\mathrm{\Sigma }_𝐪^\mathrm{L}(\mathrm{i}\omega _n)`$ $`=`$ $`{\displaystyle \underset{𝐩_1,𝐩_20}{}}|M_{𝐪,𝐩_2;𝐩_1}|^2{\displaystyle \frac{f^0(E_{𝐩_1})f^0(E_{𝐩_2})}{\mathrm{i}\omega _n(E_{𝐩_1}E_{𝐩_2})}},`$ (57)
$`\mathrm{\Sigma }_𝐪^\mathrm{B}(\mathrm{i}\omega _n)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{𝐩_1,𝐩_20}{}}|M_{𝐩_1,𝐩_2;𝐪}|^2{\displaystyle \frac{1+f^0(E_{𝐩_1})+f^0(E_{𝐩_2})}{\mathrm{i}\omega _nE_{𝐩_1}E_{𝐩_2}}}.`$ (58)
The damping of excitations is given by the imaginary part of the self-energy
$`\mathrm{\Gamma }_𝐪=\mathrm{Im}\mathrm{\Sigma }_𝐪(E_𝐪+\mathrm{i}\delta ),`$ (59)
where $`\delta +0`$. This gives the damping of Bogoliubov excitations in an optical lattice,
$`\mathrm{\Gamma }_𝐪`$ $`=`$ $`\mathrm{\Gamma }_𝐪^\mathrm{L}+\mathrm{\Gamma }_𝐪^\mathrm{B},`$ (60)
$`\mathrm{\Gamma }_𝐪^\mathrm{L}`$ $`=`$ $`\pi {\displaystyle \underset{𝐩_1,𝐩_20}{}}|M_{𝐪,𝐩_2;𝐩_1}|^2[f^0\left(E_{𝐩_2}\right)f^0\left(E_{𝐩_1}\right)]\delta \left(E_𝐪E_{𝐩_1}+E_{𝐩_2}\right),`$ (61)
$`\mathrm{\Gamma }_𝐪^\mathrm{B}`$ $`=`$ $`{\displaystyle \frac{\pi }{2}}{\displaystyle \underset{𝐩_1,𝐩_20}{}}|M_{𝐩_1,𝐩_2;𝐪}|^2[1+f^0(E_{𝐩_1})+f^0(E_{𝐩_2})]\delta \left(E_𝐪E_{𝐩_1}E_{𝐩_2}\right).`$ (62)
Landau damping $`\mathrm{\Gamma }_𝐪^\mathrm{L}`$ is expected to be dominant at higher temperatures where there are thermally excited quasiparticles. In contrast, Beliaev damping $`\mathrm{\Gamma }_𝐪^\mathrm{B}`$ is due to a decay process and can arise in the absence of thermally-excited excitations. Beliaev damping is possible even at $`T=0`$ \[i.e. $`f^0(E_𝐩)=0`$\] and is expected to be dominant at low temperatures.
### V.2 Landau damping in a 1D optical lattice
In this subsection, we calculate the Landau damping given by Eq. (61) for a 1D optical lattice. For this, the energy conservation condition $`E_𝐪+E_𝐩=E_{𝐪+𝐩+𝐆}=E_{𝐪+𝐩}`$ needs to be satisfied, where $`𝐆`$ is a reciprocal lattice vector. The solution of the energy conservation condition $`E_q+E_p=E_{q+p}`$ for a 1D optical lattice is illustrated in Fig. 5 Peierls . First, we draw an excitation spectrum and mark $`q`$ on the line. Then we draw the spectrum line again with $`q`$ as origin. If those two lines intersect, solutions consistent with energy conservation condition for three excitation process exist. The intersection can be taken as $`(q+p,E_q+E_p=E_{q+p})`$. Clearly, the condition for these two dispersion curves to have intersections requires that the dispersion relation $`E_q`$ first bends up as $`q`$ increases, before bending over. From Fig. 1, we see that the 1D optical lattice dispersion relation $`E_q`$ has this feature only for $`\alpha <6`$ and thus Landau damping can occur only when $`\alpha <6`$. If the intersection $`(q+p,E_q+E_p)`$ is outside of the first Brillouin zone, it has to be reduced in the first Brillouin zone by subtracting a reciprocal lattice vector $`G_n=\frac{2\pi n}{d}`$ ($`n`$ is an integer), corresponding to an Umklapp process Peierls .
The values of $`(q,p)`$ satisfying the energy conservation condition for 1D optical lattice are shown in Fig. 6. For a given value of $`q`$, we see that as $`\alpha 6`$, the value of $`p`$ decreases to zero. There is no solution for $`\alpha >6`$, indicating the disappearance of the damping for an excitation $`E_q`$ with any value of $`q`$. In Fig. 5, the curves of the solution of the energy conservation condition never cross the dashed line at $`q+p=\pi /d`$. Thus Umklapp scattering processes ($`G_n0`$) do not contribute to Landau damping in our present single-band model.
When the excitation of wavevector $`q`$ has a wavelength much larger than the thermal excitation $`p`$ (i.e., when $`qp`$, $`\pi /d`$), one finds that $`\frac{\mathrm{d}E_p}{\mathrm{d}p}=c`$ from the energy conservation condition $`E_q+E_p=E_{q+p}`$ in Eq. (61). The Landau damping of the excitation $`E_q=cq`$ comes from absorbing a thermal excitation $`E_p`$ with a group velocity equal to $`c`$ Pitaevskii&Levinson . This condition requires
$`{\displaystyle \frac{\mathrm{d}E_p}{\mathrm{d}p}}`$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{E_p}}{E_p}}4Jd\mathrm{sin}{\displaystyle \frac{pd}{2}}\mathrm{cos}{\displaystyle \frac{pd}{2}}`$ (63)
$`=`$ $`c\sqrt{2Jd^2Un^{\mathrm{c0}}}.`$
This equation can be written in the form
$`16X^2+8(\alpha 2)X+(\alpha ^26\alpha )=0,`$ (64)
where $`X\mathrm{sin}^2\frac{pd}{2}`$. The unique value of $`p_0`$ such that $`\frac{\mathrm{d}E_{p_0}}{\mathrm{d}p_0}`$ in Eq. (63) equals $`c`$ is found to be given by the condition
$`\mathrm{sin}^2\left({\displaystyle \frac{p_0d}{2}}\right)={\displaystyle \frac{(\alpha 2)+\sqrt{2(\alpha +2)}}{4}}.`$ (65)
This expression is only valid for $`\alpha `$ smaller than 6, such that $`p_0q`$.
For $`qp`$, where we can use the following approximations,
$`u_q`$ $``$ $`\sqrt{{\displaystyle \frac{Un^{\mathrm{c0}}}{2cp}}}\left(1+{\displaystyle \frac{cp}{2Un^{\mathrm{c0}}}}\right),`$ (66)
$`v_q`$ $``$ $`\sqrt{{\displaystyle \frac{Un^{\mathrm{c0}}}{2cp}}}\left(1{\displaystyle \frac{cp}{2Un^{\mathrm{c0}}}}\right),`$ (67)
$`u_{q+p}`$ $``$ $`u_p{\displaystyle \frac{J\mathrm{sin}(pd)}{2u_p}}{\displaystyle \frac{(Un^{\mathrm{c0}})^2}{E_p^3}}qd,`$ (68)
$`v_{q+p}`$ $``$ $`v_p{\displaystyle \frac{J\mathrm{sin}(pd)}{2v_p}}{\displaystyle \frac{(Un^{\mathrm{c0}})^2}{E_p^3}}qd,`$ (69)
the matrix element in Eq. (61) reduces to
$`M_{q,p;q+p}`$ $`=`$ $`2U\sqrt{{\displaystyle \frac{n^{\mathrm{c0}}}{I}}}\left[\left(u_pu_{q+p}+v_pv_{q+p}v_pu_{q+p}\right)u_q\left(u_pu_{q+p}+v_pv_{q+p}u_pv_{q+p}\right)v_q\right]`$ (70)
$``$ $`2U\sqrt{{\displaystyle \frac{n^{\mathrm{c0}}}{I}}}\left({\displaystyle \frac{ϵ_p^0}{E_p}}+{\displaystyle \frac{Un^{\mathrm{c0}}}{2E_p}}{\displaystyle \frac{(Un^{\mathrm{c0}})^2}{E_p^2}}{\displaystyle \frac{Jd}{c}}\mathrm{sin}pd\right)\sqrt{{\displaystyle \frac{cq}{2Un^{\mathrm{c0}}}}}.`$
From the energy conservation condition
$`c={\displaystyle \frac{\mathrm{d}E_p}{\mathrm{d}p}}={\displaystyle \frac{\stackrel{~}{E}_p}{E_p}}2Jd\mathrm{sin}pd,`$ (71)
the matrix element can be approximated by
$`M_{q,p;q+p}=U\left({\displaystyle \frac{ϵ_p^0}{E_p}}+{\displaystyle \frac{E_p}{\stackrel{~}{E}_p}}\right)\left({\displaystyle \frac{q}{2m^{}c}}\right)^{1/2}{\displaystyle \frac{N_{\mathrm{c0}}^{1/2}}{I}},`$ (72)
where $`N_{\mathrm{c0}}n^{\mathrm{c0}}I`$. This is the analogue of Eq. (10.76) in Ref. Pethick&Smith . Using the delta function for energy conservation to integrate over $`p`$, the Landau damping is finally reduced to
$`\mathrm{\Gamma }_q^\mathrm{L}={\displaystyle \frac{\alpha }{8(\alpha +2)}}\left({\displaystyle \frac{E_{p_0}}{ϵ_{p_0}^0}}\right)^3\left({\displaystyle \frac{ϵ_{p_0}^0}{E_{p_0}}}+{\displaystyle \frac{E_{p_0}}{\stackrel{~}{E}_{p_0}}}\right)^2{\displaystyle \frac{\beta JUqd}{\mathrm{sinh}^2\frac{\beta E_{p_0}}{2}}}.`$ (73)
One finds there is no damping when $`\alpha >6`$ (see Fig. 3 for $`\alpha `$ as a function of $`s`$ and $`T`$). In fact, $`\mathrm{\Gamma }_q^\mathrm{L}`$ diverges at $`\alpha =6`$, but this is due to the 1D nature of our system. This divergence does not occur in 2D and 3D lattices.
### V.3 Landau damping in 2D and 3D optical lattices
We next discuss the energy conservation condition in a 2D optical lattice. We imagine a surface in a three dimensional space which satisfies the Bogoliubov dispersion relation $`(q_x,q_y,E_𝐪)`$. Then, we draw a new dispersion $`E_𝐪`$ in the three dimensional space, with a point on the surface $`(q_{1x},q_{1y},E_{𝐪_1})`$ as origin. That is, we draw $`(q_x,q_y,E_{𝐪𝐪_1}+E_{𝐪_1})`$. If those two surfaces intersect, the energy conservation condition $`E_{𝐪_1}+E_{𝐪_2}=E_{𝐪_1+𝐪_2}`$ is satisfied, the intersection being given by $`(q_{1x}+q_{2x},q_{1x}+q_{2y},E_{𝐪_1+𝐪_2})`$. For this condition to be satisfied, the surface $`(q_x,q_y,E_𝐪)`$ has to be above the other surface $`(q_x,q_y,E_{𝐪𝐪_1}+E_𝐪)`$ around $`(q_{1x},q_{1y},E_{𝐪_1})`$. Since the Bogoliubov spectrum is phonon like $`E_𝐪cq`$ for small $`q`$, the maximum gradient of $`E_𝐪`$ at $`(q_{1x},q_{1y},E_{𝐪_1})`$ must be greater than $`c`$. Since $`|_𝐪E_𝐪|_{𝐪=𝐪_1}`$ is the maximum gradient of $`E_𝐪`$ at $`𝐪_1`$, this condition is equivalent to the requirement:
$`|_𝐪E_𝐪|_{𝐪=𝐪_1}=2Jd{\displaystyle \frac{\stackrel{~}{E_{𝐪_1}}}{E_{𝐪_1}}}\sqrt{\mathrm{sin}^2q_{1x}d+\mathrm{sin}^2q_{1y}d}>c.`$ (74)
Equation (74) is the 2D version of the condition for the Bogoliubov spectrum of a 1D optical lattice to have anomalous dispersion. If Eq. (74) is satisfied, the energy conservation condition $`E_{𝐪_1}+E_{𝐪_2}=E_{𝐪_1+𝐪_2}`$ can be satisfied. An excitation $`E_{𝐪_1}`$ can then decay into $`E_{𝐪_1+𝐪_2}`$ by absorbing $`E_{𝐪_2}`$ (Landau), or an excitation $`E_{𝐪_1+𝐪_2}`$ can decay into two excitations $`E_{𝐪_1}`$ and $`E_{𝐪_2}`$ (Beliaev).
The condition for such a $`𝐪_1`$ to exist is that the maximum value of $`|_𝐪E_𝐪|`$ as a function of $`(q_x,q_y)`$ is greater than $`c`$. Due to $`|_𝐪E_𝐪|_{𝐪=0}=c`$, we only need to consider the condition that $`|_𝐪E_𝐪|`$ takes its maximum at $`𝐪0`$. When $`|_𝐪E_𝐪|`$ has its maximum value,
$`_𝐪|_𝐪E_𝐪|`$ $`=`$ $`{\displaystyle \frac{2Jd^2\mathrm{sin}q_xd}{\sqrt{\mathrm{sin}^2q_xd+\mathrm{sin}^2q_yd}}}`$ (79)
$`\times ({\displaystyle \frac{\stackrel{~}{E}_𝐪}{E_𝐪}}\left(\begin{array}{c}\mathrm{cos}q_xd\\ \mathrm{cos}q_yd\end{array}\right){\displaystyle \frac{2JU^2(n^{\mathrm{c0}})^2}{E_𝐪^3}}(\mathrm{sin}^2q_xd+\mathrm{sin}^2q_yd)\left(\begin{array}{c}1\\ 1\end{array}\right))=0.`$
From Eq. (79), $`|_𝐪E_𝐪|`$ has its maximum value when $`\mathrm{cos}q_xd=\mathrm{cos}q_yd`$, i.e., at the two values $`q_x=\pm q_y`$. If we assume $`q_x=q_y`$ and define $`u\mathrm{sin}^2q_xd`$, Eq. (79) reduces to
$`u=\mathrm{sin}^2q_xd={\displaystyle \frac{(3\alpha 4)+\sqrt{5\alpha ^2+24\alpha +16}}{16}}.`$ (80)
From Eq. (80), one sees that $`u`$ decreases as $`\alpha `$ increases, vanishing when $`\alpha =12`$. Therefore, we conclude that the Landau damping when $`\alpha 12`$ is due to excitations with momentum $`q_x=\pm q_y`$, and all damping processes in a 2D optical lattice will vanish when $`\alpha >12`$. We have confirmed this analytical result by numerically solving the energy conservation condition.
The condition for the disappearance of Landau damping in a 3D optical lattice can be derived by generalizing the procedure described above for a 2D optical lattice. One sees that $`|_𝐪E_𝐪|`$ has its maximum when $`\mathrm{cos}q_xd=\mathrm{cos}q_yd=\mathrm{cos}q_zd`$, i.e., $`q_x=\pm q_y=\pm q_z`$ and
$`\mathrm{sin}^2q_xd={\displaystyle \frac{3(\alpha 2)+\sqrt{5\alpha ^2+36\alpha +36}}{24}}.`$ (81)
One finds that $`u0`$ when $`\alpha 18`$, and damping in a 3D optical lattice vanishes when $`\alpha 18`$. As in the 2D case, Landau damping when $`\alpha 18`$ only occurs for an excitation with momentum $`q_x=\pm q_y=\pm q_z`$.
We discuss the energy conservation condition in a 2D optical lattice in detail. We only consider the damping of long wavelength phonon $`𝐪`$. Using the approximation for the long wavelength phonon $`E_𝐪cq`$ and $`E_{𝐪+𝐩}E_𝐪+_𝐩E_𝐩𝐪`$, the energy conservation condition $`E_𝐪+E_𝐪=E_{𝐪+𝐩}`$ can be written as
$`{\displaystyle \frac{\alpha }{2}}\left(q_x^2+q_y^2\right)={\displaystyle \frac{\stackrel{~}{E}_𝐩^2}{E_𝐩^2}}\left(\mathrm{sin}^2(p_xd)q_x^2+2\mathrm{sin}(p_xd)\mathrm{sin}(p_yd)q_xq_y+\mathrm{sin}^2(p_yd)q_y^2\right).`$ (82)
When $`q_x>0`$ and $`q_y=0`$, Eq. (82) can be solved easily. Defining $`X\mathrm{sin}^2\frac{p_xd}{2}`$ and $`Y\mathrm{sin}^2\frac{p_yd}{2}`$, the solution of Eq. (82) is
$`Y=\left(X+{\displaystyle \frac{\alpha }{4}}\right)+{\displaystyle \frac{1}{4}}\sqrt{{\displaystyle \frac{\alpha ^3}{8X^28X+\alpha }}}.`$ (83)
One can confirm that Eq. (83) reduces to the 1D result in Eq. (65) when $`Y=0`$.
Equation (83) is plotted in Fig. 7 for several values of $`\alpha `$. We see that as $`\alpha 6`$ the line in the $`(p_x,p_y)`$ plane which satisfies the energy conservation condition shrinks and vanishes when $`\alpha >6`$. Therefore, the Landau damping of an excitation with $`q_y=0`$ disappears when $`\alpha >6`$.
$`p_0`$ in Fig. 2.8 is given by Eq. (65).
For a long wavelength phonon $`𝐪`$ with $`q_x=q_y>0`$, we solve the energy conservation condition numerically. The solution is shown in Fig. 8. There is no solution when $`\alpha >12`$ as expected. Figures 7 and 8 clearly show that the threshold value of $`\alpha `$ for the disappearance of damping strongly depends on the direction of $`𝐪`$ due to the anisotropy of 2D square lattice. This result also holds for 3D square lattice.
As for the 1D case, $`\mathrm{\Gamma }_𝐪^\mathrm{L}`$ becomes larger than $`E_𝐪`$ around the threshold value of $`\alpha `$ in 2D and 3D optical lattices. In this case, we cannot use the simple Golden Rule expression of the Landau damping given by Eq. (61). We have to extend it using higher order perturbation theory in order to calculate the Landau damping when $`\alpha `$ is close to the threshold value Leggett .
### V.4 Beliaev damping in a 1D optical lattice
We briefly discuss the Beliaev damping of the Bloch-Bogoliubov excitation in this subsection. The Beliaev damping is due to spontaneous decay of an excitation into two excitations, and thus we need to satify the energy conservation condition $`E_𝐪=E_{𝐪𝐩}+E_𝐩`$. We focus on a 1D case in the following.
In Fig. 9, the solution of the energy conservation condition for 1D optical lattice $`E_q=E_{qp}+E_p`$ is shown in a $`(q,p)`$ plane. As predicted in Sec. V.2, one finds that the curve of the solution shrinks as $`\alpha `$ increases and vanishes when $`\alpha >6`$, which indicates the disappearance of the Beliaev damping for $`\alpha >6`$.
For a fixed value of $`\alpha <6`$, Beliaev damping is only possible when $`q`$ is between the threshold momenta $`q_0`$ and $`q_\mathrm{c}`$ shown in Fig. (9). At the threshold momenta $`q_0`$ and $`q_\mathrm{c}`$, two excitations $`E_q`$ and $`E_{qp}`$ created by the decay of an excitation $`E_q`$ have the same velocity. This effect was first predicted by Pitaevskii in 1959 for the phonon excitation in superfluid $`{}_{}{}^{4}\mathrm{He}`$ Pitaevskii . The two created excitations have the same quasi-momentum $`q_\mathrm{c}/2`$ and energy $`E_{q_\mathrm{c}/2}`$ at $`q=q_\mathrm{c}`$. At $`q=q_0`$, one of the generated excitations is a phonon having the sound velocity $`c`$. Therefore, the other one also has the velocity equal to the sound velocity $`c`$.
In addition to Landau damping and Beliaev damping, one also has intercollisional damping arising from two body collisions which transfer atoms between the condensate and thermal cloud at finite temperatures WG ; ZNG . Such processes also involve the energy conservation condition for three-excitation processes. Thus, the intercollisional damping also disappears when $`\alpha 6D`$ in a $`D`$-dimensional optical lattice.
## VI conclusion
In conclusion, we have given a detailed treatment of the damping of Bogoliubov excitations associated with Bose condensates in an optical lattice at finite temperature using the tight-binding Bose-Hubbard model, extending the work we reported in Ref. TsuchiyaGriffin .
We have used the Popov approximation in the Bose-Hubbard model to extend the usual theory to finite temperatures. As a by-product, we have calculated the number of condensate atoms per lattice site as a function of both the temperature and the lattice well depth $`s`$. These results may be of general interest.
We calculated the damping of Bloch-Bogoliubov excitations in 1D, 2D and 3D optical lattices. All previous work Smerzi3 ; Wu2 on this problem only considered dynamical instabilities. These studies did not include the dissipation processes we have considered. We find that the Bogoliubov-Popov excitation spectrum $`E_𝐪`$ must exhibit “anomalous dispersion” for damping processes to occur. This is analogous to the case of phonon damping in superfluid $`{}_{}{}^{4}\mathrm{He}`$. In the absense of this “bending-up” of the low $`𝐪`$ spectrum, energy conservation cannot be satisfied. As a consequence, excitation damping is absent when $`\alpha =Un^{\mathrm{c0}}/J>6D`$ where $`D`$ is the dimension of the optical lattice.
The first studies Inguscio2 of damping of excitations were limited to 1D optical lattices along the axis of a cigar-shape magnetic trap, but one would need a much tighter magnetic trap (in the radial direction) for our 1D model results to apply. In the more recent experiments by Stöferle et. al. Esslinger2 , a 3D optical lattice is prepared first, and the lattice potential depths in two lattice axes are then made much larger than the third one to produce 2D array of tightly bound 1D optical lattices. An analogous 2D optical lattice can be formed by choosing a much larger lattice potential depth along one lattice axis than that in other two axes. This effectively 2D optical lattice might be better for checking our theoretical predictions than the 2D poriodic array of 1D tubes used in Ref. Esslinger1 , since excitations along the direction perpendicular to the 2D lattice potential can be neglected.
Due to our use of a tight-binding approximation, our results are not directly applicable to the damping of excitations found in very weak optical lattices ($`s<1`$) Porto . Extension of our calculations to such weak optical lattices would be clearly of interest thesis .
###### Acknowledgements.
S.T. would like to thank M. Smith, E. Taylor, and D. Luxat for discussions. S.T. was supported by the Japan Society for Promotion of Science (JSPS). A.G. was supported by funds from NSERC of Canada.
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# Exact soliton solution of Spin Chain with a external magnetic field in linear wave background
## I Introduction
The concept of soliton in spin chain which exhibits both coherent and chaotic structures depending on the nature of the magnetic interactionsTrullinger ; Kivshar ; Ablowitz ; Fokas has received considerable attention for decades. Soliton in quasi one-dimensional magnetic systems is no longer a theoretical concept but can be probed by neutron inelastic scattering Kjems78 ; Boucher90 , Mossbauer linewidth measurements Thiel81 , and electron spin resonance Asano00 . One of the integrable models for spin chain is Landau-Lifshitz equation landau which has been studied by a variety of techniques such as inverse scattering transformationAblowitz ; Takhtajan ; Chen , Darboux transformation Huang , Riemann-Hilbert approachliu1 , etc. Exact soliton solutions are reported for the isotropic Takhtajan ; Tjon ; Fogedby ; Laksmanan ; Shimizu and anisotropic Chen ; Huang ; Mikeska ; Li ; liu1 spin chains as well. The other integrable model is a type of Nonlinear-Schrődinger equation studied in Ref. Ablowitz . It should be noted that all these solutions are obtained in the ground state background. The exact soliton solutions in linear wave background corresponding to magnon in quantum theoryMajlis has received no attention yet. The main goal of this paper is to search for new solutions of spin chain driven by external magnetic field in linear wave background. We obtain exact solution of N-soliton train in terms of a simple, straightforward Darboux transformationMatveev ; C.H. Gu . Particularly the time-evolution of one and two-soliton is analyzed in terms of the general solution.
The outline of this paper is organized as follows: In Sec. II the Darboux transformation is explained in detail and the general N-soliton solution is obtained. In Sec. III we discuss a special case. Exact one-soliton solution in linear wave background is obtained. The dispersion law and soliton energy are investigated as well. Sec. IV is devoted to general two-soliton solution and soliton collisions. Finally, our concluding remarks are given in Sec. V.
## II Exact solution of N-soliton train
Our starting Hamiltonian describing the anisotropic spin chain with a external magnetic field can be written as
$`\widehat{H}`$ $`=g\mu _BB{\displaystyle \underset{j}{}}\widehat{S}_j^zJ^{}{\displaystyle \underset{j,\epsilon }{}}\widehat{S}_j^z\widehat{S}_{j+\epsilon }^z`$
$`{\displaystyle \frac{1}{2}}J{\displaystyle \underset{j,\epsilon }{}}\left(\widehat{S}_j^+\widehat{S}_{j+\epsilon }^{}+\widehat{S}_j^{}\widehat{S}_{j+\epsilon }^+\right),`$ (1)
where $`\widehat{S}_j(\widehat{S}_j^x,\widehat{S}_j^y,\widehat{S}_j^z)`$ with $`j=1,2,\mathrm{},N`$ are spin operators, $`J^{}>J`$ $`>0`$ is the pair interaction parameter, $`g`$ the Lande factor and $`\mu _B`$ is the Bohr magneton, $`B`$ is the external magnetic field. With the help of Holstein-Primokoff Holstein transformation $`\widehat{S}_j^z=Sa_j^{}a_j`$, $`\widehat{S}_j^+\sqrt{2S}\left(1\frac{1}{4S}a_j^{}a_j\right)a_j`$, $`\widehat{S}_j^{}\sqrt{2S}a_j^{}\left(1\frac{1}{4S}a_j^{}a_j\right)`$, the Hamiltonian in Eq. (1) reduces to
$`\widehat{H}`$ $`=2J^{}S^2Ng\mu _BBSN+g\mu _BB{\displaystyle \underset{j}{}}a_j^{}a_j`$
$`+S{\displaystyle \underset{j,\epsilon }{}}\left\{J^{}(a_j^{}a_j+a_{j+\epsilon }^{}a_{j+\epsilon })J(a_ja_{j+\epsilon }^{}+a_j^{}a_{j+\epsilon })\right\}`$
$`J^{}{\displaystyle \underset{j,\epsilon }{}}a_j^{}a_ja_{j+\epsilon }^{}a_{j+\epsilon }+{\displaystyle \frac{1}{4}}J{\displaystyle \underset{j,\epsilon }{}}[a_j^{}a_{j+\epsilon }^{}a_{j+\epsilon }a_{j+\epsilon }`$
$`+a_{j+\epsilon }^{}a_j^{}a_ja_j+a_j^{}a_j^{}a_ja_{j+\epsilon }+a_{j+\epsilon }^{}a_{j+\epsilon }^{}a_{j+\epsilon }a_j]`$
where the boson operators $`a_j`$ are assumed to satisfy the usual commutation relation $`[a_j,a_j^{}^{}]=\delta _{jj^{}}`$, $`[a_j^{},a_j^{}^{}]=[a_j,a_j^{}]=0`$. The equation of motion for the operator $`a_j`$ on the nth site is $`i\mathrm{}\frac{}{t}a_j=[a_j,\widehat{H}]`$. At low temperatures, the operator $`a_j`$ can be treated as a classical vector such that $`a_ja(x.t)`$. So that the equation of motion in a continuum spin chain under a magnetic field can be obtained as a Nonlinear-Schrődinger type:
$$i\frac{}{t}a=\beta _0\frac{^2}{x^2}a+2\beta _0\beta _1^2a\left|a\right|^2+2\beta _2a,$$
(2)
where
$$\beta _0=\frac{2JS}{\mathrm{}},\beta _1=\sqrt{\frac{J^{^{}}J}{JS}},\beta _2=\frac{g\mu _BB+4\left(J^{^{}}J\right)S}{2\mathrm{}},$$
here $`J^{^{}}>J>0`$(easy-axis). In this paper we will present a systematic method to construct general expressions of one- and two-soliton solutions embedded in a linear wave background for Eq. (2) and their novel properties.
By employing Ablowitz-Kaup-Newell-Segur technique one can construct the linear eigenvalue problem for Eq. (2) as follows
$$\psi _x=U\psi ,\text{ }\psi _t=F\psi ,$$
(3)
where $`\psi =\left(\begin{array}{c}\psi _1\hfill \\ \psi _2\hfill \end{array}\right)`$, $`U`$ and $`F`$ can be given in the forms
$`U`$ $`=\lambda \sigma _3+P,`$
$`F`$ $`=i\left(2\lambda ^2\beta _0+\beta _2\right)\sigma _3+i2\lambda \beta _0Pi\beta _0\left[P^2+P_x\right]\sigma _3,`$
with
$$\sigma _3=\left(\begin{array}{cc}1\hfill & \text{ }0\hfill \\ 0\hfill & 1\hfill \end{array}\right),P(x,t)=\left(\begin{array}{cc}0\hfill & \beta _1a\hfill \\ \beta _1\overline{a}\hfill & 0\hfill \end{array}\right),$$
and the overbar denotes the complex conjugate. Thus Eq. (2) can be recovered from the compatibility condition $`U_tF_x+[U,F]=0`$. Based on the Lax pair (3), we can obtain general one- and two-soliton solution embedded in a linear wave background by using a straightforward Darboux transformationMatveev ; C.H. Gu .
Consider the following transformation
$$\mathrm{\Psi }=\left(\lambda IK\right)\psi ,\text{ }K=H\mathrm{\Lambda }H^1,\text{ }\mathrm{\Lambda }=\text{diag}(\lambda _1,\lambda _2),$$
(4)
where $`H`$ is a nonsingular matrix which satisfies
$$H_x=\sigma _3H\mathrm{\Lambda }+PH.$$
(5)
Letting
$$\mathrm{\Psi }_x=U_1\mathrm{\Psi },$$
(6)
where $`U_1=\lambda \sigma _3+P_1`$, $`P_1=\left(\begin{array}{cc}0\hfill & \beta _1a_1\hfill \\ \beta _1\overline{a}_1\hfill & \text{ }0\hfill \end{array}\right)`$, and with the help of Eqs. (3), (4) and (5), we obtain the Darboux transformation for Eq. (2) from Eq. (6) in the form
$$P_1=P+[\sigma _3,K]\text{.}$$
(7)
It is easy to verify that, if $`\psi =\left(\begin{array}{c}\psi _1\hfill \\ \psi _2\hfill \end{array}\right)`$ is a eigenfunction of Eq. (3) with eigenvalue $`\lambda =\lambda _1`$, then $`\left(\begin{array}{c}\overline{\psi }_2\hfill \\ \overline{\psi }_1\hfill \end{array}\right)`$ is also the eigenfunction, however with eigenvalue $`\overline{\lambda }_1`$. Thus if taking
$$H=\left(\begin{array}{cc}\psi _1\hfill & \overline{\psi }_2\hfill \\ \psi _2\hfill & \text{ }\overline{\psi }_1\hfill \end{array}\right),\mathrm{\Lambda }=\left(\begin{array}{cc}\lambda _1\hfill & \text{ }0\hfill \\ 0\hfill & \overline{\lambda }_1\hfill \end{array}\right),$$
(8)
which ensures that Eq. (5) is held, we can obtain
$$K_{sl}=\overline{\lambda }_1\delta _{sl}+\left(\lambda _1+\overline{\lambda }_1\right)\frac{\psi _s\overline{\psi }_l}{\psi ^T\overline{\psi }},\text{ }s,l=1,2,$$
(9)
where $`\psi ^T\overline{\psi }=\left|\psi _1\right|^2+\left|\psi _2\right|^2,`$ and Eq. (7) becomes
$$a_1=a+\frac{2}{\beta _1}\left(\lambda _1+\overline{\lambda }_1\right)\frac{\psi _1\overline{\psi }_2}{\psi ^T\overline{\psi }},$$
(10)
where $`\psi =(\psi _1,\psi _2)^T`$ is the eigenfunction of Eq. (3) corresponding to the eigenvalue $`\lambda _1`$ for the solution $`a`$. Thus by solving the Eq. (3) which is a first-order linear differential equation, we can generate a new solution $`a_1`$ of the Eq. (2) from a known solution $`a`$ which is usually called “seed” solution.
To obtain exact $`N`$-order solution of Eq. (2), we firstly rewrite the Darboux transformation in Eq. (10) as in the form
$$a_1=a+\frac{2}{\beta _1}\left(\lambda _1+\overline{\lambda }_1\right)\frac{\psi _1[1,\lambda _1]\overline{\psi }_2[1,\lambda _1]}{\psi [1,\lambda _1]^T\overline{\psi }[1,\lambda _1]},$$
(11)
where $`\psi [1,\lambda ]=(\psi _1[1,\lambda ],\psi _2[1,\lambda ])^T`$ denotes the eigenfunction of Eq. (3) corresponding to eigenvalue $`\lambda `$. Then repeating above the procedure for $`N`$ times, we can obtain the exact $`N`$-order solution
$$a_N=a+\frac{2}{\beta _1}\underset{m=1}{\overset{N}{}}(\lambda _m+\overline{\lambda }_m)\frac{\psi _1[m,\lambda _m]\overline{\psi }_2[m,\lambda _m]}{\psi [m,\lambda _m]^T\overline{\psi }[m,\lambda _m]},$$
(12)
where
$`\psi [m,\lambda ]`$ $`=\left(\lambda K\left[m1\right]\right)\mathrm{}\left(\lambda K\left[1\right]\right)\psi [1,\lambda ],`$
$`K_{sl}\left[j^{}\right]`$ $`=\overline{\lambda }_j^{}\delta _{sl}+\left(\lambda _j^{}+\overline{\lambda }_j^{}\right){\displaystyle \frac{\psi _s[j^{},\lambda _j^{}]\overline{\psi }_l[j^{},\lambda _j^{}]}{\psi [j^{},\lambda _j^{}]^T\overline{\psi }[j^{},\lambda _j^{}]}},`$
here $`\psi [j^{},\lambda ]`$ is the eigenfunction corresponding to $`\lambda _j^{}`$ for $`a_{j^{}1}`$ with $`a_0a`$ and $`s,l=1,2`$, $`j^{}=1,2,\mathrm{},m1`$, $`m=2,3,\mathrm{},N`$. Thus if choosing a “seed” as the basic initial solution, by solving linear characteristic equation system (3), one can construct a set of new solutions for Eq. (2) by employing the formula (12).
As an example, we give the exact expression of one- and two-soliton solutions in a linear wave background for Eq. (2) respectively and analyze its properties. For this propose, we take the initial “seed” as $`a=A_ce^{ik_cxi\omega _ct}`$ which is a linear-wave solution and satisfies the nonlinear dispersion relation $`\omega _c=\beta _0\left(k_c^22\beta _1^2A_c^2\right)2\beta _2`$, where $`A_c`$ and $`\omega _c`$ are the arbitrary real constants. Substituting the initial seed into (3) and solving the linear equation, by tedious calculation we obtain the expression of eigenfunction corresponding to eigenvalue $`\lambda `$ in the form
$`\psi _1`$ $`=\beta _1A_cC_1\mathrm{exp}\mathrm{\Theta }_1+L_1C_2\mathrm{exp}\mathrm{\Theta }_2,`$
$`\psi _2`$ $`=L_1C_1\mathrm{exp}\left(\mathrm{\Theta }_2\right)\beta _1A_cC_2\mathrm{exp}\left(\mathrm{\Theta }_1\right),`$ (13)
where
$`\mathrm{\Theta }_1`$ $`={\displaystyle \frac{1}{2}}i\phi +iM\left(x\gamma t\right),\text{ }\mathrm{\Theta }_2={\displaystyle \frac{1}{2}}i\phi iM\left(x\gamma t\right),`$
$`L`$ $`=\lambda i{\displaystyle \frac{k_c}{2}}iM,\text{ }M={\displaystyle \frac{1}{2}}\sqrt{\left(k_c+i2\lambda \right)^2+4\beta _1^2A_c^2},`$
$`\gamma `$ $`=\left(k_ci2\lambda \right)\beta _0,\text{ }\phi =k_cx\omega _ct.`$
here the parameters $`C_1`$ and $`C_2`$ are the arbitrary complex constants. Following the expression of eigenfunction and with the aid of the formulas (12), we can obtain exactly the one- and two-soliton solutions embedded in a linear wave background, respectively.
## III One-soliton solution embedded in linear wave background
Taking the spectral parameter $`\lambda =\lambda _1\frac{A_{s,1}}{2}+i\frac{k_{s,1}}{2}`$ in Eq. (13) and substituting them into Eq. (11), we obtain the one-soliton solution in the linear wave background as follows
$`a_1`$ $`=e^{i\phi }\{A_c+{\displaystyle \frac{A_{s,1}}{\beta _1\mathrm{\Delta }_1}}[\beta _1^2A_c^2e^{i\mathrm{\Phi }_1}+|L_1|^2e^{i\mathrm{\Phi }_1}`$
$`2\beta _1A_c(\mathrm{Re}L_1\mathrm{cosh}\theta _1+i\mathrm{Im}L_1\mathrm{sinh}\theta _1)]\}`$ (14)
where
$`\mathrm{\Delta }_1`$ $`=\left(\left|L_1\right|^2+\beta _1^2A_c^2\right)\mathrm{cosh}\theta _12\beta _1A_c\left(\mathrm{Re}L_1\right)\mathrm{cos}\mathrm{\Phi },`$
$`\theta _1`$ $`=2\mathrm{Im}M_1\left(xV_1t\right)x_0,`$
$`\mathrm{\Phi }_1`$ $`=2\mathrm{Re}M_1\left(xV_2t\right)\phi _0,`$
$`V_1`$ $`={\displaystyle \frac{\mathrm{Im}\left(M_1\gamma _1\right)}{\mathrm{Im}M_1}},\text{ }V_2={\displaystyle \frac{\mathrm{Re}\left(M_1\gamma _1\right)}{\mathrm{Re}M_1}},\text{ }\phi =k_cx\omega _ct,`$
with the parameters $`M_1=\frac{1}{2}\sqrt{(k_c+i2\lambda _1)^2+4\beta _1^2A_c^2}`$, $`L_1=\lambda _1ik_c/2iM_1`$, and $`\gamma _1=\left(k_ci2\lambda _1\right)\beta _0`$. The parameters $`x_0`$ and $`\phi _0`$ represent the initial center position and initial phase, which are determined by $`x_0=\mathrm{ln}\left|C_2/C_1\right|`$, and $`\phi _0=\mathrm{arg}\left(C_2/C_1\right),`$ respectively, where $`C_1,C_2`$ are the arbitrary complex constants. It is worth to note that because $`x_0`$ and $`\phi _0`$ are determined by the value $`C_2/C_1,`$ they in fact depend on only one arbitrary complex parameter. Not loss of generality, we take $`C_1=1`$.
The exact solution $`a_1`$ in Eq. (14) describes a soliton solution of anisotropic spin chain embedded in a linear-wave background with the soliton amplitude $`\frac{A_{s,1}}{\beta _1}`$, the width $`\frac{1}{2\mathrm{Im}M_1}`$, the wavenumber $`k_1=2\mathrm{Re}M_1`$, the frequency $`\mathrm{\Omega }_1=2\mathrm{Re}\left(M_1\gamma _1\right)`$, and the envelope velocity $`V_1=\mathrm{Im}\left(M_1\gamma _1\right)/\mathrm{Im}M_1`$. As the linear-wave amplitude vanishes, namely $`A_c=0`$, this solution in Eq. (14) reduces to the solution in the from
$$a_{1\text{sol}}=\frac{A_se^{i\left[k_{s,1}x\mathrm{\Omega }_{s,1}t+\phi _0\right]}}{\beta _1\mathrm{cosh}\left[A_{s,1}\left(xV_{s,1}tx_0^{}\right)\right]},$$
(15)
where $`x_0^{}`$ is determined by $`x_0^{}=\frac{1}{A_{s,1}}\mathrm{ln}\left|C_2\right|`$. The solution $`a_{1\text{sol}}`$ in Eq. (15) describes a bright soliton solution with maximal amplitude $`\frac{A_{s,1}}{\beta _1}`$, the width $`\frac{1}{A_{s,1}}`$, envelope velocity $`V_{s,1}=2\beta _0k_{s,1}`$, and center position $`x_0^{}`$. The frequency $`\mathrm{\Omega }_{s,1}=\frac{1}{\mathrm{}}[g\mu _BB+4\left(J^{^{}}J\right)S2JSA_{s,1}^2+\frac{1}{2}\frac{\mathrm{}^2}{4JS}V_{s,1}^2]`$ and wavenumber $`k_{s,1}=\frac{V_{s,1}}{2\beta _0}`$ of the “carrier wave” are related by the dispersion law $`\mathrm{\Omega }_{s,1}=\beta _0(k_{s,1}^2A_{s,1}^2)2\beta _2=\beta _0(k_{s,1}^2A_{s,1}^2)+\frac{1}{\mathrm{}}[g\mu _BB+4\left(J^{^{}}J\right)S]`$. It is shown that the external magnetic field $`B`$ change the frequency but not the amplitude. Then the soliton energy is seen to be
$$E_1=\mathrm{}\mathrm{\Omega }_{s,1}=g\mu _BB+4\left(J^{^{}}J\right)S2JSA_{s,1}^2+\frac{1}{2}m^{}V_{s,1}^2$$
where an effective mass $`m^{}`$ of soliton is $`\frac{\mathrm{}^2}{4JS}`$. We also notice that the velocity of the “carrier wave”, that is the phase velocity of the soliton $`\frac{\mathrm{\Omega }_{s,1}}{k_{s,1}}=\frac{V_{s,1}}{2}\frac{\beta _0A_{s,1}^2+2\beta _2}{k_{s,1}}`$, has a negative correction $`\frac{\beta _0A_{s,1}^2+2\beta _2}{k_{s,1}}`$ for the half of envelope velocity, whereas the soliton group velocity $`\frac{d\mathrm{\Omega }_{s,1}}{dk_{s,1}}=V_{s,1}`$ coincides with the envelope velocity. On the other hand when the soliton amplitude vanishes, namely $`A_{s,1}=0`$, the solution $`a_1`$ in Eq. (14) reduces to the linear-wave solution $`a=A_ce^{i\phi },`$ where $`\phi =k_cx\omega _ct`$ and group velocity $`V_c=\frac{d\omega _c}{dk_c}=2\beta _0k_c`$ coming from the nonlinear dispersion relation $`\omega _c=\beta _0\left(k_c^22\beta _1^2A_c^2\right)2\beta _2`$.
From the expression of $`M_1`$ we can directly see that $`M_1`$ is the pure real number when $`k_c=k_{s,1}`$ and $`A_{s,1}^2<4\beta _1^2A_c^2`$. The condition $`A_{s,1}^2<4\beta _1^2A_c^2=\frac{4\eta (J^{^{}}J)}{JS}A_c^2`$ is a stability criterion which is related to anisotropic parameter $`(J^{^{}}J)`$ and the amplitude $`A_c`$ of the linear wave. It is worth to point out that the condition $`k_c=k_{s,1}`$ implies the equal group velocities for both soliton and linear wave.
## IV Two-soliton solution embedded in linear wave background
According to the general formalism in section II it is easy to construct the two-soliton solution in the linear wave background as follows:
$$a_2=a_1+e^{i\phi }\frac{A_{s,2}}{\beta _1\mathrm{cosh}\mathrm{\Gamma }}e^{i\mathrm{arg}h_2},$$
(16)
where
$`h_2`$ $`={\displaystyle \frac{\left(\lambda _2+\overline{\lambda }_1\right)\mathrm{exp}\left(i\phi \right)\left(\lambda _1+\overline{\lambda }_1\right)\rho _1+\left(\lambda _2\lambda _1\right)\left|\rho _1\right|^2\rho _2}{\left(\lambda _2\lambda _1\right)+\left(\lambda _2+\overline{\lambda }_1\right)\left|\rho _1\right|^2+\left(\lambda _1+\overline{\lambda }_1\right)\overline{\rho }_1\rho _2}},`$
$`\mathrm{\Gamma }`$ $`=\mathrm{ln}\left|h_2\right|,\text{ }\rho _n={\displaystyle \frac{L_n\beta _1A_ce^{\theta _ni\mathrm{\Phi }_n}}{\beta _1A_c+L_ne^{\theta _ni\mathrm{\Phi }_n}}},`$
with the notations
$`\theta _n`$ $`=2\mathrm{Im}M_n\left(xV_{1,n}t\right)x_{0,n},`$
$`\mathrm{\Phi }_n`$ $`=2\mathrm{Re}M_n\left(xV_{2,n}\right)t\phi _{0,n},`$
$`V_{1,n}`$ $`={\displaystyle \frac{\mathrm{Im}\left(M_n\gamma _n\right)}{\mathrm{Im}M_n}},V_{2,n}={\displaystyle \frac{\mathrm{Re}\left(M_n\gamma _n\right)}{\mathrm{Re}M_n}},`$
$`\phi `$ $`=k_cx\omega _ct.`$
here the parameters $`\lambda _n=\mu A_{s,n}/2+ik_{s,n}/2`$, $`L_n=\lambda _nik_c/2iM_n`$, $`M_n=\frac{1}{2}\sqrt{(k_c+i2\lambda _n)^2+4\beta _1^2A_c^2}`$ and $`\gamma _n=\left(k_ci2\lambda _n\right)\beta _0`$. The parameters $`x_{0,n}`$ and $`\phi _{0,n}`$ represent the initial center position and initial phase, which are determined by $`x_{0,n}=\mathrm{ln}\left|C_{2,n}/C_{1,n}\right|`$, and $`\phi _{0,n}=\mathrm{arg}\left(C_{2,n}/C_{1,n}\right)`$, respectively, where $`C_{1,n},C_{2,n}`$ are the arbitrary complex constants with $`n=1,2`$. With the similar reason as in the case of one-soliton solution, we often set $`C_{1,n}=1`$ and $`C_{2,n}`$ are the arbitrary complex constants.
As the linear wave amplitude vanishes $`A_c=0,`$ we have the general two-bright soliton solution from Eq. (16) in the form
$$a_{2\text{sol}}=\frac{2}{\beta _1\mathrm{\Delta }_2}(G_1e^{i[k_{s,2}x\mathrm{\Omega }_{s,2}t+\phi _{0,2}]}+G_2e^{i[k_{s,1}x\mathrm{\Omega }_{s,1}t+\phi _{0,1}]}),$$
(17)
where
$`G_1`$ $`=\left[\left(\zeta _1\mathrm{Re}\zeta _3\right)\mathrm{cosh}\theta _1+i\mathrm{Im}\zeta _3\mathrm{sinh}\theta _1\right],`$
$`G_2`$ $`=\left[\left(\zeta _2\mathrm{Re}\zeta _3\right)\mathrm{cosh}\theta _2i\mathrm{Im}\zeta _3\mathrm{sinh}\theta _2\right],`$
$`\mathrm{\Delta }_2`$ $`=\zeta _4\mathrm{cosh}\theta _1\mathrm{cosh}\theta _2`$
$`A_{s,1}A_{s,2}\left[\mathrm{cosh}\left(\theta _1+\theta _2\right)+\mathrm{cos}\left(\mathrm{\Phi }_1\mathrm{\Phi }_2\right)\right],`$
here $`\zeta _1=A_{s,2}\left|\lambda _2+\overline{\lambda }_1\right|^2`$, $`\zeta _2=A_{s,1}\left|\lambda _2+\overline{\lambda }_1\right|^2`$, $`\zeta _3=A_{s,1}A_{s,2}\left(\overline{\lambda }_2+\lambda _1\right)`$, $`\zeta _4=2\left|\lambda _2+\overline{\lambda }_1\right|^2`$. The solution $`a_{2\text{sol}}`$ in Eq. (17) is the general form of two-bright soliton solution for Eq. (2) which describes the interaction of two one-bright soliton solutions with the maximal amplitudes $`\frac{A_{s,n}}{\beta _1}`$, the width $`\frac{1}{A_{s,n}}`$, envelope velocity $`V_{s,n}=2\beta _0k_{s,n}`$, $`n=1,2`$, respectively. The frequency $`\mathrm{\Omega }_{s,n}=\frac{1}{\mathrm{}}[g\mu _BB+4\left(J^{^{}}J\right)S2JSA_{s,n}^2+\frac{1}{2}\frac{\mathrm{}^2}{4JS}V_{s,n}^2]`$ and wavenumber $`k_{s,n}=\frac{V_{s,n}}{2\beta _0}`$ of each “carrier wave” are related by the dispersion law $`\mathrm{\Omega }_{s,n}=\beta _0\left(k_{s,n}^2A_{s,n}^2\right)2\beta _2`$. Then the energy of each soliton is seen to be
$$E_n=\mathrm{}\mathrm{\Omega }_{s,n}=g\mu _BB+4\left(J^{^{}}J\right)S2JSA_{s,n}^2+\frac{1}{2}m^{}V_{s,n}^2,$$
with $`n=1,2`$, where $`m^{}=\frac{\mathrm{}^2}{4JS}`$ denotes the effective mass of soliton. We also notice that the velocity of each “carrier wave”, that is the phase velocity of each soliton $`\frac{\mathrm{\Omega }_{s,n}}{k_{s,n}}=\frac{V_{s,n}}{2}\frac{\beta _0A_{s,n}^2+2\beta _2}{k_{s,n}}`$, has a negative correction $`\frac{\beta _0A_{s,n}^2+2\beta _2}{k_{s,n}}`$ for the half of envelope velocity, whereas the group velocity of each soliton $`\frac{d\mathrm{\Omega }_{s,n}}{dk_{s,n}}=V_{s,n}`$ coincides with the envelope velocity of each soliton. When the amplitude $`A_{s,n}=0,`$ $`n=1,2,`$ the solution $`a_2`$ in Eq. (16) reduces to the linear wave solution. Therefore, in general, the solution $`a_2`$ in Eq. (16) represents the interaction of two one-soliton solution in a linear wave background.
As the discussion for Eq. (14), we consider the case of $`k_c=k_{s,1}=k_{s,2}`$. Hence we have $`M_n=\frac{1}{2}\sqrt{4\beta _1^2A_c^2A_{s,n}^2}`$, $`n=1,2`$. From this expression, we are easy to see that the condition $`A_{s,n}^2<4\beta _1^2A_c^2=\frac{4\eta \left(J^{^{}}J\right)}{JS}A_c^2`$, $`n=1,2`$, becomes a stability criterion which is related to anisotropic parameter $`\left(J^{^{}}J\right)`$ and the amplitude $`A_c`$ of the linear wave. It is also worth to see that the condition $`k_c=k_{s,1}=k_{s,2}`$ implies the equal group velocities both for each soliton and linear wave.
## V Conclusion
In this paper we obtain exact N-soliton solution for anisotropic spin chain driven by a external magnetic field in linear wave background in terms of a simple, straightforward Darboux transformation. As a special case the explicit one- and two-soliton solution dressed by the linear wave corresponding to magnon in quantum theory is obtained analytically and its property is discussed in detail. The frequency $`\mathrm{\Omega }_{s,n}`$, wavenumber $`k_{s,n}`$, and the dispersion law of each “carrier wave” are also studied. We obtain explicitly the energy $`E_n`$ of each soliton and the effective mass $`m^{}`$ of soliton. Our result show that the stability criterion of soliton is related with anisotropic parameter and the amplitude of the linear wave.
## VI Acknowledgment
This work was support by the NSF of China under Grant No. 10075032.
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# 𝒩={1,2} supersymmetric vacua of IIA supergravity and 𝑆𝑈(2) structures
## 1 Introduction and summary
String theory compactifications in the presence of fluxes possess a number of phenomenologically attractive features, a fact which has led to their intensive study in recent years. String theory is approximated in the low-energy effective-theory limit by supergravity and in many situations of physical interest it has proven extremely fruitful to study the properties of supersymmetric supergravity solutions. The presence of some amount of unbroken supersymmetry above a certain low (compared to the Planck mass) energy scale is desirable, if one wishes to avoid issues of stability.
The subject of supersymmetric supergravity compactifications with fluxes is a particularly old one. Nevertheless, it has only recently been realized (starting with ; see for a recent review and references) that the machinery of $`G`$-structures can be a powerful tool in classifying and constructing supergravity solutions. From this point of view, a $`G`$-structure is the natural generalization of the concept of special holonomy to the case where nontrivial fluxes, i.e. nonzero vevs of the antisymmetric tensor fields, are present.
In we presented a classification of $`𝒩=1`$ supersymmetric solutions of IIA supergravity of the form of a warped product $`AdS_4\times _\omega X_6`$, where $`X_6`$ is a six-dimensional compact manifold of $`SU(3)`$ structure, generalizing the work of . The manifold $`X_6`$ was constrained to be ‘half-flat’ of a certain type. For further related work on IIA compactifications from the point of view of the four-dimensional effective field theory see . Type IIA compactifications have also been considered in the context of $`G`$-structures in . The recent paper analyzes type II $`𝒩=1`$ supersymmetry using the concept of generalized $`G`$-structures –we will come back to this in the next paragraph. Supersymmetric $`AdS_4`$ solutions are of additional interest as they are expected to be dual to certain three-dimensional superconformal field theories .
As will be explained in the following, for the backgrounds considered here $`𝒩=1`$ supersymmetry implies that the Majorana-Weyl supersymmetry parameter $`ϵ`$ is of the form
$`ϵ`$ $`=\theta _+(\alpha \eta _{1+}+\delta \eta _2)+\mathrm{c}.\mathrm{c}.,`$ (1)
where $`\alpha `$, $`\delta `$, are functions on $`X_6`$, $`\eta _{1+}`$ ($`\eta _2`$) is a globally-defined, chiral (antichiral) unimodular spinor (and therefore nowhere-vanishing) on $`X_6`$ and $`\theta `$ is a Killing spinor of $`M_{1,3}`$<sup>1</sup><sup>1</sup>1$`\eta _{1,2+}`$ are related to $`\eta _{1,2}`$ by complex conjugation. Equation (1) represents one linear combination of $`\eta _{1,2}`$ and corresponds to $`𝒩=1`$ in $`d=4`$. . The existence of $`\eta _1`$ implies that the structure group of $`X_6`$ is reduced to $`SU(3)`$. If in addition $`\eta _{1,2}`$ are nowhere-parallel, the structure group is further reduced to $`SU(2)`$. Relaxing the condition that $`\eta _{1,2}`$ should be linearly-independent everywhere on $`X_6`$ would lead to a situation which can be thought of as a so-called ‘generalized $`SU(3)`$ structure’ on $`X_6`$ : at generic points in $`X_6`$ the two $`SU(3)`$ structures associated with each of the two internal spinors have a common subgroup, which defines an $`SU(2)`$ structure on $`X_6`$. However, at the points where the two spinors become parallel the structure group collapses to $`SU(3)`$.
Supersymmetric vacua on manifolds of $`SU(2)`$ structure restrict the choice for the fluxes similar to their $`SU(3)`$-structure counterparts. In addition, the requirement of an $`SU(2)`$ structure imposes a strong constraint on the internal manifold<sup>2</sup><sup>2</sup>2A necessary and sufficient condition for the structure group of a manifold $`X_6`$ of $`SU(3)`$ structure to further reduce to $`SU(2)`$, is the existence of a globally-defined nowhere-vanishing vector field on $`X_6`$. This is equivalent to the requirement that the Euler characteristic vanish, $`\chi (X_6)=0`$. and one may hope that a classification can proceed much more explicitly in this case. However the situation is much more difficult to analyze in practice, and this subject has received much less attention in the literature, because of the multitude of flux-components which arise in decomposing the supergravity fields in terms of irreducible $`SU(2)`$ representations.
In the present paper we examine $`𝒩=1`$ supersymmetric type IIA vacua in the case where $`X_6`$ is a compact manifold of $`SU(2)`$ structure. In we considered the case $`\eta _1=\eta _2`$. Here we will consider the other ‘extreme’ case where $`\eta _{1,2}`$ are everywhere orthogonal<sup>3</sup><sup>3</sup>3This is more restrictive than requiring that $`\eta _{1,2}`$ in (1) be nowhere-parallel.. In addition, we look for solutions with nonzero Romans’ mass. We reformulate the supersymmetry conditions in terms of $`SU(2)`$ structures in section 5. In search for explicit solutions we make some further simplifying assumptions; namely, we set all nonscalar (in the sense of irreducible $`SU(2)`$ representations) fluxes to zero and we take the dilaton to be constant. This is a consistent truncation which, however, turns out to be too stringent: as we will see there do not exist any supersymmetric vacua of this type.
A related $`𝒩=1`$ type IIA vacuum was constructed in . In that case the manifold $`X_6`$ was taken to be conformally $`R^6`$ and therefore noncompact, allowing for non-constant harmonic functions. Taking $`X_6`$ to be $`T^6`$ instead in the solution of , would imply that the warp factors and the dilaton are constant and that the Romans’ mass and all the fluxes are zero. The solution would therefore degenerate to $`R^{1,3}\times T^6`$, in agreement with our conclusion above. Type IIA supersymmetric compactifications to Minkowski space on manifolds of $`SU(2)`$ structure have also been considered in . However, all ten-dimensional vacua in that paper arise upon reduction of eleven-dimensional supergravity solutions and are therefore unrelated to the present work<sup>4</sup><sup>4</sup>4Recall that for nonzero mass parameter, as is the case here, Romans’ supergravity has no Poincaré-invariant lift to eleven dimensions; see for a recent discussion..
In section 6 we proceed to examine the case of $`𝒩=2`$ supersymmetric (warped) $`AdS_4`$ vacua. Rather than considering the most general spinor Ansatz, we will take the two Majorana-Weyl supersymmetry parameters $`ϵ_{1,2}`$ to be of the form
$`ϵ_1`$ $`=\theta _+(\alpha \eta _{1+}+\beta \eta _1)+\mathrm{c}.\mathrm{c}.`$
$`ϵ_2`$ $`=\theta _+(\gamma \eta _{2+}+\delta \eta _2)+\mathrm{c}.\mathrm{c}.,`$ (2)
where $`\alpha `$, $`\beta `$, $`\gamma `$, $`\delta `$, are functions on $`X_6`$ and $`\eta _{1,2}`$ are globally-defined, unimodular spinors on $`X_6`$. In addition we take $`\eta _{1,2}`$ to be orthogonal to each other. Requiring $`M_{1,3}`$ to be $`AdS_4`$ implies $`\alpha =\beta `$, $`\gamma =\delta `$. As we show in section 5, this requirement is again too stringent and in fact there do not exist any $`𝒩=2`$ IIA vacua of this type.
Note that for an admissible vacuum, except for the supersymmetry conditions we also require the supergravity equations of motion and the Bianchi identities for the forms to be satisfied. More generally, the no-go theorems of this paper could be by-passed by introducing additional sources, for example orientifolds, which would modify the equations-of-motion. Alternatively one may consider singular and/or noncompact ‘internal’ manifolds, higher-order stringy corrections, etc. We also emphasize that, due to the technical complexity of the task, we have not been able to analyze the most general spinor Ansatz leading to $`SU(2)`$ structures. We hope to report on that in the future.
The outline of the remainder of the paper is as follows: In the next section we review some useful facts about (Romans’) IIA supergravity. Our Ansatz for the ten-dimensional $`𝒩=1`$ background and the corresponding reduction of the supersymmetry conditions are considered in section 3. Section 4 contains a brief review of $`SU(3)`$ and $`SU(2)`$ structures in six dimensions. The main analysis of our Ansatz for $`𝒩=1`$ supersymmetric vacua is contained in section 5. Section 6 contains the analysis of our $`𝒩=2`$ Ansatz. Most of the technical details and some further useful formulæ are relegated to the four appendices.
## 2 Massive IIA
For the sake of completeness, in this section we note some known facts about Romans’ ten-dimensional supergravity. Our notation and conventions are as in to which the reader is referred for further details.
The equations of motion for the bosonic fields of massive IIA supergravity are
$`0`$ $`=R_{MN}{\displaystyle \frac{1}{2}}_M\varphi _N\varphi {\displaystyle \frac{1}{12}}e^{\varphi /2}G_{MPQR}G_N{}_{}{}^{PQR}+{\displaystyle \frac{1}{128}}e^{\varphi /2}g_{MN}G^2`$
$`{\displaystyle \frac{1}{4}}e^\varphi H_{MPQ}H_N{}_{}{}^{PQ}+{\displaystyle \frac{1}{48}}e^\varphi g_{MN}H^2`$
$`2m^2e^{3\varphi /2}B_{MP}^{}B_N^{}{}_{}{}^{P}+{\displaystyle \frac{m^2}{8}}e^{3\varphi /2}g_{MN}(B^{})^2{\displaystyle \frac{m^2}{4}}e^{5\varphi /2}g_{MN}`$ (3)
$`0`$ $`=^2\varphi {\displaystyle \frac{1}{96}}e^{\varphi /2}G^2+{\displaystyle \frac{1}{12}}e^\varphi H^2{\displaystyle \frac{3m^2}{2}}e^{3\varphi /2}(B^{})^25m^2e^{5\varphi /2}`$ (4)
$`0`$ $`=d(e^\varphi H){\displaystyle \frac{1}{2}}GG+2me^{\varphi /2}B^{}G+4m^2e^{3\varphi /2}B^{}`$ (5)
$`0`$ $`=d(e^{\varphi /2}G)HG.`$ (6)
In addition, the forms obey the Bianchi identities
$`dB^{}`$ $`=H`$
$`dH`$ $`=0`$
$`dG`$ $`=2mB^{}H.`$ (7)
To make contact with the massless IIA supergravity of one introduces a Stückelberg gauge potential $`A`$, with field strength $`F=dA`$, so that
$`mB^{}`$ $`=mB+{\displaystyle \frac{1}{2}}F.`$ (8)
In the massless limit, $`m0`$, we have $`mB^{}\frac{1}{2}F`$.
### Supersymmetry
The gravitino and dilatino supersymmetry variations read
$`\delta \mathrm{\Psi }_M=𝒟_Mϵ`$ (9)
and
$`\delta \lambda =\{{\displaystyle \frac{1}{2}}\mathrm{\Gamma }^M_M\varphi `$ $`{\displaystyle \frac{5me^{5\varphi /4}}{4}}+{\displaystyle \frac{3me^{3\varphi /4}}{8}}B_{MN}^{}\mathrm{\Gamma }^{MN}\mathrm{\Gamma }_{11}`$
$`+{\displaystyle \frac{e^{\varphi /2}}{24}}H_{MNP}\mathrm{\Gamma }^{MNP}\mathrm{\Gamma }_{11}{\displaystyle \frac{e^{\varphi /4}}{192}}G_{MNPQ}\mathrm{\Gamma }^{MNPQ}\}ϵ,`$ (10)
where
$`𝒟_M`$ $`:=\{_M{\displaystyle \frac{me^{5\varphi /4}}{16}}\mathrm{\Gamma }_M{\displaystyle \frac{me^{3\varphi /4}}{32}}B_{NP}^{}(\mathrm{\Gamma }_M{}_{}{}^{NP}14\delta _M{}_{}{}^{N}\mathrm{\Gamma }_{}^{P})\mathrm{\Gamma }_{11}`$
$`+{\displaystyle \frac{e^{\varphi /2}}{96}}H_{NPQ}(\mathrm{\Gamma }_M{}_{}{}^{NPQ}9\delta _M{}_{}{}^{N}\mathrm{\Gamma }_{}^{PQ})\mathrm{\Gamma }_{11}+{\displaystyle \frac{e^{\varphi /4}}{256}}G_{NPQR}(\mathrm{\Gamma }_M{}_{}{}^{NPQR}{\displaystyle \frac{20}{3}}\delta _M{}_{}{}^{N}\mathrm{\Gamma }_{}^{PQR})\}.`$ (11)
One can transform to the string frame by rescaling $`e_A{}_{}{}^{M}e^{\varphi /4}e_A^M`$.
### Integrability
It was shown in that imposing supersymmetry together with the equations of motion for the forms implies the dilaton equation and the Einstein equation $`E_{MN}=0`$, provided $`E_{M0}=0`$ for $`M0`$ <sup>5</sup><sup>5</sup>5Similar integrability conditions were derived in in the context of eleven-dimensional supergravity. See also for a recent general discussion..
## 3 $`𝒩=1`$ $`M_{1,3}\times _\omega X_6`$ backgrounds
Let us now assume that spacetime is of the form of a warped product $`M_{1,3}\times _\omega X_6`$, where $`M_{1,3}`$ is Minkowski or $`AdS_4`$ and $`X_6`$ is a compact manifold. The ten dimensional metric reads
$$g_{MN}(x,y)=\left(\begin{array}{cc}\mathrm{\Delta }^2(y)\widehat{g}_{\mu \nu }(x)& 0\\ 0& \rho _{mn}(y)\end{array}\right),$$
(12)
where $`x`$ is a coordinate on $`M_{1,3}`$ and $`y`$ is a coordinate on $`X_6`$. We will also assume that the forms have nonzero $`y`$-dependent components along the internal directions, except for the four-form which will be allowed to have an additional component proportional to the volume of $`M_{1,3}`$
$$G_{\mu \nu \kappa \lambda }=\sqrt{g_4}f(y)\epsilon _{\mu \nu \kappa \lambda },$$
(13)
where $`f`$ is a real scalar function on $`X_6`$. Note that with these assumptions the $`E_{M0}=0`$ for $`M0`$ condition is satisfied, and therefore we need only check supersymmetry the Bianchi identities and the equations of motion for the forms.
### 3.1 Massive $`𝒩=1`$ vacua and $`SU(2)`$ structure
On $`M_{1,3}`$ there is a pair of Weyl spinors (related by complex conjugation), each of which satisfies the Killing equation
$$\widehat{}_\mu \theta _+=W\widehat{\gamma }_\mu \theta _{};\widehat{}_\mu \theta _{}=W^{}\widehat{\gamma }_\mu \theta _+,$$
(14)
where hatted quantities are computed using the metric $`\widehat{g}_{\mu \nu }`$, and the complex constant $`W`$ is related to the scalar curvature $`\widehat{R}`$ of $`M_{1,3}`$ through $`\widehat{R}=24|W|^2`$. The reader is referred to for further details on our spinor conventions in four, six, ten dimensions.
It can been shown (see for example for a recent discussion) that the requirement of $`𝒩=1`$ supersymmetry<sup>6</sup><sup>6</sup>6This corresponds to four real supercharges; in the present paper we are counting supersymmetries according to four-dimensional conventions. implies that the Majorana-Weyl supersymmetry parameter $`ϵ`$ is decomposed under $`Spin(1,9)Spin(1,3)\times Spin(6)`$ as
$`ϵ`$ $`=\alpha (y)\theta _+\eta _{1+}+\delta (y)\theta _+\eta _2+\mathrm{c}.\mathrm{c}.,`$ (15)
where $`\alpha `$, $`\delta `$, are complex functions on $`X_6`$ and $`\eta _{1,2}`$ is a pair of globally-defined, nowhere-vanishing Weyl spinors on $`X_6`$. Moreover, $`\eta _{1,2}`$ are related to $`\eta _{1,2+}`$ by complex conjugation. Without loss of generality, we can choose $`\eta _{1,2}`$ to be of unit norm. In keeping with four-dimensional supersymmetry nomenclature, we take $`\theta `$ ($`\eta _{1,2}`$) to be anticommuting (commuting).
There are three cases according to the relation between $`\eta _1`$ and $`\eta _2`$: a) $`\eta _2`$ is everywhere parallel to $`\eta _1`$ and $`X_6`$ is of $`SU(3)`$ structure, b) $`\eta _2`$ is nowhere parallel to $`\eta _1`$ and $`X_6`$ is of $`SU(2)`$ structure, or c) at generic points in $`X_6`$ $`\eta _{1,2}`$ are linearly independent, but there exist points where $`\eta _2`$ becomes parallel to $`\eta _1`$. The latter case imposes no additional topological requirement on $`X_6`$ other than that its structure group should reduce to $`SU(3)`$.
In the present paper we will take $`\eta _2`$ to be everywhere orthogonal to $`\eta _1`$: $`\eta _{2+}^+\eta _{1+}=0`$. This is a special sub-case of b) above and therefore $`X_6`$ must be a manifold of $`SU(2)`$ structure. As we will see later in section 5, requiring in addition that the Romans’ mass be nonzero implies that up to a choice of phase which can be absorbed in the normalization of the spinors $`\eta _{1,2}`$,
$`|\alpha |=\alpha =\delta .`$ (16)
### 3.2 Reduction of the supersymmetry conditions
Substituting the spinor Ansatz (15) in the supersymmetry transformations we obtain
$`0`$ $`=\alpha _m\eta _{1+}+_m\alpha \eta _{1+}+\alpha {\displaystyle \frac{e^{\varphi /2}}{96}}H_{npq}(\gamma _m{}_{}{}^{npq}9\delta _m{}_{}{}^{n}\gamma _{}^{pq})\eta _{1+}\delta {\displaystyle \frac{me^{5\varphi /4}}{16}}\gamma _m\eta _2`$
$`+3i\delta f{\displaystyle \frac{e^{\varphi /4}}{32}}\gamma _m\eta _2+\delta {\displaystyle \frac{me^{3\varphi /4}}{32}}B_{np}^{}(\gamma _m{}_{}{}^{np}14\delta _m{}_{}{}^{n}\gamma _{}^{p})\eta _2`$
$`+\delta {\displaystyle \frac{e^{\varphi /4}}{256}}G_{npqr}(\gamma _m{}_{}{}^{npqr}{\displaystyle \frac{20}{3}}\delta _m{}_{}{}^{n}\gamma _{}^{pqr})\eta _2`$ (17)
$`0`$ $`=\delta ^{}_m\eta _{2+}+_m\delta ^{}\eta _{2+}\delta ^{}{\displaystyle \frac{e^{\varphi /2}}{96}}H_{npq}(\gamma _m{}_{}{}^{npq}9\delta _m{}_{}{}^{n}\gamma _{}^{pq})\eta _{2+}+\alpha ^{}{\displaystyle \frac{me^{5\varphi /4}}{16}}\gamma _m\eta _1`$
$`+3i\alpha ^{}f{\displaystyle \frac{e^{\varphi /4}}{32}}\gamma _m\eta _1+\alpha ^{}{\displaystyle \frac{me^{3\varphi /4}}{32}}B_{np}^{}(\gamma _m{}_{}{}^{np}14\delta _m{}_{}{}^{n}\gamma _{}^{p})\eta _1`$
$`\alpha ^{}{\displaystyle \frac{e^{\varphi /4}}{256}}G_{npqr}(\gamma _m{}_{}{}^{npqr}{\displaystyle \frac{20}{3}}\delta _m{}_{}{}^{n}\gamma _{}^{pqr})\eta _1,`$ (18)
from the ‘internal’ components of the gravitino variation and
$`0`$ $`=\alpha \mathrm{\Delta }^1W\eta _{1+}+\delta ^{}{\displaystyle \frac{me^{5\varphi /4}}{16}}\eta _{2+}5i\delta ^{}f{\displaystyle \frac{e^{\varphi /4}}{32}}\eta _{2+}\delta ^{}{\displaystyle \frac{me^{3\varphi /4}}{32}}B_{mn}^{}\gamma ^{mn}\eta _{2+}`$
$`+\alpha ^{}{\displaystyle \frac{e^{\varphi /2}}{96}}H_{mnp}\gamma ^{mnp}\eta _1\delta ^{}{\displaystyle \frac{e^{\varphi /4}}{256}}G_{mnpq}\gamma ^{mnpq}\eta _{2+}{\displaystyle \frac{1}{2}}\alpha ^{}_m(\mathrm{ln}\mathrm{\Delta })\gamma ^m\eta _1`$ (19)
$`0`$ $`=\delta ^{}\mathrm{\Delta }^1W^{}\eta _{2+}+\alpha {\displaystyle \frac{me^{5\varphi /4}}{16}}\eta _{1+}+5i\alpha f{\displaystyle \frac{e^{\varphi /4}}{32}}\eta _{1+}+\alpha {\displaystyle \frac{me^{3\varphi /4}}{32}}B_{mn}^{}\gamma ^{mn}\eta _{1+}`$
$`+\delta {\displaystyle \frac{e^{\varphi /2}}{96}}H_{mnp}\gamma ^{mnp}\eta _2\alpha {\displaystyle \frac{e^{\varphi /4}}{256}}G_{mnpq}\gamma ^{mnpq}\eta _{1+}+{\displaystyle \frac{1}{2}}\delta _m(\mathrm{ln}\mathrm{\Delta })\gamma ^m\eta _2,`$ (20)
from the noncompact piece. Note that these equations are complex. Similarly from the dilatino we obtain
$`0`$ $`={\displaystyle \frac{1}{2}}\alpha ^{}_m\varphi \gamma ^m\eta _1\alpha ^{}{\displaystyle \frac{e^{\varphi /2}}{24}}H_{mnp}\gamma ^{mnp}\eta _1\delta ^{}{\displaystyle \frac{5me^{5\varphi /4}}{4}}\eta _{2+}`$
$`+i\delta ^{}f{\displaystyle \frac{e^{\varphi /4}}{8}}\eta _{2+}\delta ^{}{\displaystyle \frac{3me^{3\varphi /4}}{8}}B_{mn}^{}\gamma ^{mn}\eta _{2+}\delta ^{}{\displaystyle \frac{e^{\varphi /4}}{192}}G_{mnpq}\gamma ^{mnpq}\eta _{2+}`$ (21)
$`0`$ $`={\displaystyle \frac{1}{2}}\delta _m\varphi \gamma ^m\eta _2+\delta {\displaystyle \frac{e^{\varphi /2}}{24}}H_{mnp}\gamma ^{mnp}\eta _2+\alpha {\displaystyle \frac{5me^{5\varphi /4}}{4}}\eta _{1+}`$
$`+i\alpha f{\displaystyle \frac{e^{\varphi /4}}{8}}\eta _{1+}\alpha {\displaystyle \frac{3me^{3\varphi /4}}{8}}B_{mn}^{}\gamma ^{mn}\eta _{1+}+\alpha {\displaystyle \frac{e^{\varphi /4}}{192}}G_{mnpq}\gamma ^{mnpq}\eta _{1+}.`$ (22)
## 4 $`SU(2)`$ reduction
The analysis of the conditions for a supersymmetric vacuum and the characterization of the solutions is greatly facilitated by using the machinery of $`G`$-structures . The existence of two globally-defined nowhere-vanishing orthogonal spinors $`\eta _{1,2}`$, as is the case here, implies the reduction of the structure group of $`X_6`$ to $`SU(2)`$. This allows us to decompose all tensors on $`X_6`$ in terms of irreducible $`SU(2)`$ representations. In the following two subsections we review some of the relevant facts about $`SU(3)`$ and $`SU(2)`$ structures, before we turn to the analysis of the conditions for an $`𝒩=1`$ supersymmetric vacuum in section 5. The details of the $`SU(2)`$ decomposition of the antisymmetric forms of IIA supergravity and the $`SU(2)`$ decomposition of the supersymmetry conditions are given in appendices C, D respectively.
### 4.1 $`SU(3)`$ structure
The existence of a nowhere-vanishing globally-defined spinor $`\eta _1`$ allows us to define the bilinears
$$J_{mn}:=i\eta _1^+\gamma _{mn}\eta _1=i\eta _{1+}^+\gamma _{mn}\eta _{1+}$$
(23)
$$\mathrm{\Omega }_{mnp}:=\eta _1^+\gamma _{mnp}\eta _{1+};\mathrm{\Omega }_{mnp}^{}=\eta _{1+}^+\gamma _{mnp}\eta _1.$$
(24)
Note that $`J_{mn}`$ thus defined is real and $`\mathrm{\Omega }`$ ($`\mathrm{\Omega }^{}`$) is imaginary (anti-) self-dual.
$$\mathrm{\Omega }_{mnp}=\frac{i}{6}\sqrt{\rho _6}\epsilon _{mnpijk}\mathrm{\Omega }^{ijk}.$$
(25)
We choose to normalize
$$\eta _{1+}^+\eta _{1+}=\eta _1^+\eta _1=1.$$
(26)
Using (94) one can prove that $`J`$, $`\mathrm{\Omega }`$ satisfy
$$J_m{}_{}{}^{n}J_{n}^{}{}_{}{}^{p}=\delta _m^p$$
(27)
$$(\mathrm{\Pi }^+)_m{}_{}{}^{n}\mathrm{\Omega }_{npq}^{}=\mathrm{\Omega }_{mpq};(\mathrm{\Pi }^{})_m{}_{}{}^{n}\mathrm{\Omega }_{npq}^{}=0,$$
(28)
where
$$(\mathrm{\Pi }^\pm )_m{}_{}{}^{n}:=\frac{1}{2}(\delta _m{}_{}{}^{n}iJ_m{}_{}{}^{n})$$
(29)
are the projection operators onto the holomorphic/antiholomorphic parts. In other words, $`J`$ defines an almost complex structure with respect to which $`\mathrm{\Omega }`$ is $`(3,0)`$. Moreover (using (94) again) it follows that
$`\mathrm{\Omega }J`$ $`=0`$
$`\mathrm{\Omega }\mathrm{\Omega }^{}`$ $`={\displaystyle \frac{4i}{3}}J^3.`$ (30)
Therefore $`J`$, $`\mathrm{\Omega }`$, completely specify an $`SU(3)`$ structure on $`X_6`$.
Further useful relations can be found in .
### 4.2 $`SU(2)`$ structure
The existence of two orthogonal unimodular globally-defined spinors $`\eta _1`$, $`\eta _2`$ on $`X_6`$ allows us to define two distinct $`SU(3)`$ structures
$`J_{mn}:=i\eta _1^+\gamma _{mn}\eta _1;\mathrm{\Omega }_{mnp}:=\eta _1^+\gamma _{mnp}\eta _{1+}`$ (31)
and
$`J_{mn}^{}:=i\eta _2^+\gamma _{mn}\eta _2;\mathrm{\Omega }_{mnp}^{}:=\eta _2^+\gamma _{mnp}\eta _{2+}.`$ (32)
Each of these satisfies all the properties of $`SU(3)`$ structures given in the preceding section. It can be shown, using the above definitions and the Fierz identities in appendix A, that
$`J`$ $`={\displaystyle \frac{i}{2}}KK^{}+\stackrel{~}{J}`$
$`J^{}`$ $`={\displaystyle \frac{i}{2}}KK^{}\stackrel{~}{J},`$ (33)
where
$`K_m`$ $`:=\eta _2^+\gamma _m\eta _{1+}`$ (34)
and
$`\iota _K\stackrel{~}{J}=0.`$ (35)
The complex vector $`K`$ satisfies
$`K_mK^m=0;K_m^{}K^m=2`$ (36)
and is holomorphic with respect to $`J`$,
$`(\mathrm{\Pi }^+)_m{}_{}{}^{n}K_{n}^{}=K_m;(\mathrm{\Pi }^{})_m{}_{}{}^{n}K_{n}^{}=0.`$ (37)
It follows that the two-form $`\stackrel{~}{J}`$ is $`(1,1)`$ with respect to the almost complex structure $`J`$.
The two holomorphic three-forms can be expressed in the following way
$`\mathrm{\Omega }`$ $`=iK\omega `$
$`\mathrm{\Omega }^{}`$ $`=iK\omega ^{},`$ (38)
where
$`\omega _{mn}`$ $`:=i\eta _1^+\gamma _{mn}\eta _2`$ (39)
satisfies
$`\iota _K\omega =\iota _K^{}\omega =0`$ (40)
and is holomorphic with respect to $`J`$,
$`(\mathrm{\Pi }^+)_m{}_{}{}^{n}\omega _{np}^{}=\omega _{mp};(\mathrm{\Pi }^{})_m{}_{}{}^{n}\omega _{np}^{}=0.`$ (41)
It is straightforward to show that $`\stackrel{~}{J}`$, $`\omega `$ specify an $`SU(2)`$ structure. Indeed, it follows from the above formulæ that
$`\stackrel{~}{J}`$ $`\omega =0`$
$`\omega \omega ^{}`$ $`=2\stackrel{~}{J}\stackrel{~}{J}.`$ (42)
The complex vector $`K`$ specifies an almost product structure
$`R_m{}_{}{}^{n}:=K_mK^n+K_m^{}K^n\delta _m{}_{}{}^{n},`$ (43)
such that
$`R_m{}_{}{}^{n}R_{n}^{}{}_{}{}^{p}=\delta _m{}_{}{}^{p}.`$ (44)
Further useful relations are given in appendix B.
## 5 Analysis of the conditions
To analyze the supersymmetry conditions of section 3.2 it is useful to note that equations (17, 18) can be cast in the form
$`U_m\eta _{1+}+U_{mn}\gamma ^n\eta _1=0,`$ (45)
whereas equations (19-22) can be written as
$`V\eta _{1+}+V_m\gamma ^m\eta _1=0.`$ (46)
The explicit expressions for the $`U`$’s and $`V`$’s can be read off from the expressions in appendix D. The tensors $`U`$, $`V`$ further decompose into directions parallel and perpendicular to the complex vector $`K`$ defined in 4.2. The components perpendicular to $`K`$ are further decomposed in terms of irreducible $`SU(2)`$ representations. The details of the decomposition are given in appendix C. To illustrate the procedure, let us decompose
$`V_m=v_m+vK_m,`$ (47)
where $`K^mv_m=K^mv_m=0`$. We also noted that in the decomposition of $`V_m`$ there are no terms proportional to $`K_m^{}`$, due to 100. It follows that the scalar content of (46) is equivalent to:
$`V=v=0.`$ (48)
We proceed similarly for all other representations.
Let us consider the scalar component of the supersymmetry equations first. It is straightforward to show that if there exists a point $`y_0`$ in $`X_6`$ such that $`|\alpha (y_0)||\delta (y_0)|`$, equations (19-22) imply that $`m=W=0`$. I.e. the mass parameter vanishes and the space $`M_{1,3}`$ reduces to Minkowski. We would like to look for massive solutions of IIA and hence, up to phases which can be absorbed in the normalizations of $`\eta _{1,2}`$ we can take:
$`|\alpha |=\alpha =\delta ,`$ (49)
at each point in $`X_6`$. Let us first analyze the supersymmetry equations (19\- 22), considering each irreducible $`SU(2)`$ representation in turn. The decompositions of all antisymmetric tensors in terms of irreducible $`SU(2)`$ representations can be found in appendix C. One can show that the solution is equivalent to the following conditions:
The 1
$`f`$ $`=0`$
$`W`$ $`={\displaystyle \frac{i\mathrm{\Delta }}{8}}(me^{3\varphi /4}b_2^{}+{\displaystyle \frac{i}{12}}e^{\varphi /4}g_2^{}),`$ (50)
$`mb_1`$ $`=0`$
$`me^{3\varphi /4}b_3`$ $`={\displaystyle \frac{i}{2}}(d_K\varphi _{}d_K^{}\varphi _{}),`$ (51)
where $`d_K:=K^m_m`$, $`d_K^{}:=K^m_m`$ and $`\varphi _\pm :=\varphi \pm 4\mathrm{l}\mathrm{n}\mathrm{\Delta }`$. Also
$`e^{\varphi /4}g_1`$ $`=16(me^{5\varphi /4}+{\displaystyle \frac{1}{4}}d_K\varphi _{}+{\displaystyle \frac{1}{4}}d_K^{}\varphi _{})`$
$`g_3`$ $`=0,`$ (52)
$`e^{\varphi /2}h_1=e^{\varphi /2}h_2^{}`$ $`={\displaystyle \frac{9me^{3\varphi /4}}{2}}b_2{\displaystyle \frac{ie^{\varphi /4}}{8}}g_2`$
$`e^{\varphi /2}h_3`$ $`=12i(me^{5\varphi /4}+{\displaystyle \frac{1}{8}}d_K\varphi _{}+{\displaystyle \frac{1}{4}}d_K^{}\varphi _+).`$ (53)
The 2
$`\stackrel{~}{}_m^+\mathrm{ln}\mathrm{\Delta }`$ $`={\displaystyle \frac{1}{4}}\stackrel{~}{}_m^+\varphi {\displaystyle \frac{me^{3\varphi /4}}{8}}\omega _m{}_{}{}^{n}\stackrel{~}{b}_{1n}^{}{\displaystyle \frac{e^{\varphi /4}}{64}}\stackrel{~}{g}_{2m}^{},`$ (54)
where $`\stackrel{~}{}_m`$ is defined in (146) and $`\stackrel{~}{}_m^\pm :=(\stackrel{~}{\mathrm{\Pi }}^\pm )_m{}_{}{}^{n}\stackrel{~}{}_{n}^{}`$. Moreover
$`me^{3\varphi /4}\stackrel{~}{b}_{2m}`$ $`=me^{3\varphi /4}\stackrel{~}{b}_{1m}^{}{\displaystyle \frac{e^{\varphi /4}}{32}}\omega _m^{}{}_{}{}^{n}(\stackrel{~}{g}_{1n}\stackrel{~}{g}_{2n}^{}),`$ (55)
$`e^{\varphi /2}\stackrel{~}{h}_{1m}`$ $`={\displaystyle \frac{e^{\varphi /4}}{4}}(\stackrel{~}{g}_{1m}\stackrel{~}{g}_{2m}^{})`$
$`e^{\varphi /2}\stackrel{~}{h}_{2m}`$ $`=3i\{2\omega _m{}_{}{}^{n}\stackrel{~}{}_{n}^{}\varphi +3me^{3\varphi /4}\stackrel{~}{b}_{1m}+{\displaystyle \frac{e^{\varphi /4}}{16}}\omega _m{}_{}{}^{n}({\displaystyle \frac{1}{2}}\stackrel{~}{g}_{1n}^{}\stackrel{~}{g}_{2n})\},`$ (56)
The 3
This representation drops out of equations (19-22).
Next we turn to the equations (17,18). The fact that $`\eta _{1,2}`$ are unimodular implies $`(\eta _1^+\eta _1)=0`$ and $`(\eta _2^+\eta _2)=0`$ which, taking (17,18) into account, can be seen to be equivalent to
$`\alpha =\mathrm{constant}\times \mathrm{\Delta }^{1/2}.`$ (57)
In addition, the orthogonality of $`\eta _{1,2}`$ implies $`(\eta _1^+\eta _2)=0`$ which, taking (17,18) into account, leads to the condition
$`h_1=h_2^{}={\displaystyle \frac{ie^{3\varphi /4}}{4}}g_2.`$ (58)
Comparing with (53, 50) we conclude that $`mb_2=ie^{\varphi /2}g_2/12`$ and $`W=0`$. Note that the $`\mathrm{𝟐}`$ representation drops out of the orthogonality constraints.
To summarize the conditions so far:
$`f,W`$ $`=0`$ (59)
In addition, in form notation,
$`mB^{}`$ $`=\left[{\displaystyle \frac{im}{4}}\mathrm{Im}(b_1^{})+{\displaystyle \frac{1}{64}}e^{\varphi /2}(\stackrel{~}{g}_1\stackrel{~}{g}_2^{})\right]K+\mathrm{c}.\mathrm{c}.`$
$`+m\stackrel{~}{b}{\displaystyle \frac{e^{3\varphi /4}}{4}}\stackrel{~}{J}\mathrm{Im}(d_K\varphi _{})+{\displaystyle \frac{e^{\varphi /2}}{48}}\mathrm{Im}(\omega g_2^{})`$
$`H`$ $`={\displaystyle \frac{1}{3}}\left\{\stackrel{~}{h}+{\displaystyle \frac{e^{3\varphi /4}}{16}}\mathrm{Im}(\omega g_2^{})+3ie^{\varphi /2}\stackrel{~}{J}\left[me^{5\varphi /4}+{\displaystyle \frac{1}{8}}d_K\varphi _{}+{\displaystyle \frac{1}{4}}d_K^{}\varphi _+\right]\right\}K+\mathrm{c}.\mathrm{c}.`$
$`{\displaystyle \frac{e^{3\varphi /4}}{16}}\stackrel{~}{J}\mathrm{Im}(\stackrel{~}{g}_1+\stackrel{~}{g}_2)+2iIm\left\{e^{\varphi /2}d^+\varphi +{\displaystyle \frac{3e^{5\varphi /4}}{8}}b_1^{}+{\displaystyle \frac{e^{3\varphi /4}}{32}}({\displaystyle \frac{1}{2}}\stackrel{~}{g}_1+\stackrel{~}{g}_2)\right\}KK^{}`$
$`G`$ $`=\left[me^\varphi +{\displaystyle \frac{e^{\varphi /4}}{2}}\mathrm{Re}(d_K\varphi _{})\right]\stackrel{~}{J}\stackrel{~}{J}`$
$`{\displaystyle \frac{i}{32}}\{(\stackrel{~}{g}_1\stackrel{~}{g}_2)\stackrel{~}{J}K+\mathrm{c}.\mathrm{c}.\}{\displaystyle \frac{i}{12}}[\stackrel{~}{g}+{\displaystyle \frac{1}{4}}\mathrm{Re}(\omega g_2^{})]KK^{},`$ (60)
where we have defined $`b_{1m}^{}:=\omega _m^{}{}_{}{}^{n}b_{1n}^{}`$. The differential equations (17, 18) determine the specific $`SU(2)`$ structure of $`X_6`$ and impose further constraints. In order to search for explicit solutions we will now make some simplifying assumptions. Namely, we will assume that only scalar fluxes are present (i.e. only the $`\mathrm{𝟏}`$ components of the form fields are nonzero) and that the dilaton is constant <sup>7</sup><sup>7</sup>7In the remainder of this section we will absorb all dilaton dependence by a field redefinition of the forms.. We use (17, 18) to read off the exterior derivatives on $`\mathrm{\Omega }`$, $`J`$ and $`K`$:
$`d\mathrm{\Omega }`$ $`=i(m\omega +{\displaystyle \frac{g_2}{48}}\stackrel{~}{J})KK^{}+{\displaystyle \frac{g_2}{12}}\stackrel{~}{J}\stackrel{~}{J}`$
$`dJ`$ $`={\displaystyle \frac{1}{24}}\mathrm{Re}(\omega g_2^{})\mathrm{Re}(K)`$
$`dK`$ $`=mKK^{}+{\displaystyle \frac{i}{24}}\mathrm{Re}(\omega g_2^{}).`$ (61)
Moreover, we can read off the action of the exterior derivative on the $`SU(2)`$ structure
$`d\stackrel{~}{J}`$ $`=0`$
$`d\omega `$ $`={\displaystyle \frac{g_2}{48}}\stackrel{~}{J}K^{}.`$ (62)
Combining all the above, we note that the nilpotency of the exterior derivative $`d^2=0`$ implies
$`g_2=0.`$ (63)
It is now straightforward to see that the Bianchi identities imply that all fluxes are zero and $`m=0`$, contrary to our assumption. We therefore conclude that there are no solutions obeying our simplified Ansatz.
## 6 $`𝒩=2`$ $`AdS_4`$ vacua and $`SU(2)`$ structure
In this section we will search for $`𝒩=2`$ supersymmetric vacua of the type $`AdS_4\times _\omega X_6`$. In order to simplify the computation, we will not consider the most general spinor Ansatz. Instead we demand that the background be invariant under two supersymmetries $`ϵ_{1,2}`$ of the form
$`ϵ_1`$ $`=\alpha (y)\theta _+\eta _{1+}+\beta (y)\theta _+\eta _1+\mathrm{c}.\mathrm{c}.,`$ (64)
and
$`ϵ_2`$ $`=\gamma (y)\theta _+\eta _{2+}+\delta (y)\theta _+\eta _2+\mathrm{c}.\mathrm{c}.,`$ (65)
where $`\alpha `$, $`\beta `$, $`\gamma `$, $`\delta `$, are complex functions on $`X_6`$ and $`\eta _{1,2}`$ are globally-defined, unimodular spinors on $`X_6`$. In addition we will take $`\eta _{1,2}`$ to be orthogonal to each other. Consequently, $`X_6`$ must be a manifold of $`SU(2)`$ structure. As we will see later in section 6.2, under these assumptions supersymmetry implies that up to a choice of phase which can be absorbed in the normalizations of the spinors $`\eta _{1,2}`$,
$`\alpha `$ $`=\beta `$
$`\gamma `$ $`=\delta .`$ (66)
### 6.1 Reduction of the supersymmetry conditions
Substituting the spinor Ansatz (64) in the supersymmetry transformations we obtain
$`0`$ $`=\alpha _m\eta _{1+}+_m\alpha \eta _{1+}+\alpha {\displaystyle \frac{e^{\varphi /2}}{96}}H_{npq}(\gamma _m{}_{}{}^{npq}9\delta _m{}_{}{}^{n}\gamma _{}^{pq})\eta _{1+}\beta {\displaystyle \frac{me^{5\varphi /4}}{16}}\gamma _m\eta _1`$
$`+3i\beta f{\displaystyle \frac{e^{\varphi /4}}{32}}\gamma _m\eta _1+\beta {\displaystyle \frac{me^{3\varphi /4}}{32}}B_{np}^{}(\gamma _m{}_{}{}^{np}14\delta _m{}_{}{}^{n}\gamma _{}^{p})\eta _1`$
$`+\beta {\displaystyle \frac{e^{\varphi /4}}{256}}G_{npqr}(\gamma _m{}_{}{}^{npqr}{\displaystyle \frac{20}{3}}\delta _m{}_{}{}^{n}\gamma _{}^{pqr})\eta _1`$ (67)
$`0`$ $`=\beta ^{}_m\eta _{1+}+_m\beta ^{}\eta _{1+}\beta ^{}{\displaystyle \frac{e^{\varphi /2}}{96}}H_{npq}(\gamma _m{}_{}{}^{npq}9\delta _m{}_{}{}^{n}\gamma _{}^{pq})\eta _{1+}+\alpha ^{}{\displaystyle \frac{me^{5\varphi /4}}{16}}\gamma _m\eta _1`$
$`+3i\alpha ^{}f{\displaystyle \frac{e^{\varphi /4}}{32}}\gamma _m\eta _1+\alpha ^{}{\displaystyle \frac{me^{3\varphi /4}}{32}}B_{np}^{}(\gamma _m{}_{}{}^{np}14\delta _m{}_{}{}^{n}\gamma _{}^{p})\eta _1`$
$`\alpha ^{}{\displaystyle \frac{e^{\varphi /4}}{256}}G_{npqr}(\gamma _m{}_{}{}^{npqr}{\displaystyle \frac{20}{3}}\delta _m{}_{}{}^{n}\gamma _{}^{pqr})\eta _1,`$ (68)
from the ‘internal’ components of the gravitino variation and
$`0`$ $`=\alpha \mathrm{\Delta }^1W\eta _{1+}+\beta ^{}{\displaystyle \frac{me^{5\varphi /4}}{16}}\eta _{1+}5i\beta ^{}f{\displaystyle \frac{e^{\varphi /4}}{32}}\eta _{1+}\beta ^{}{\displaystyle \frac{me^{3\varphi /4}}{32}}B_{mn}^{}\gamma ^{mn}\eta _{1+}`$
$`+\alpha ^{}{\displaystyle \frac{e^{\varphi /2}}{96}}H_{mnp}\gamma ^{mnp}\eta _1\beta ^{}{\displaystyle \frac{e^{\varphi /4}}{256}}G_{mnpq}\gamma ^{mnpq}\eta _{1+}{\displaystyle \frac{1}{2}}\alpha ^{}_m(\mathrm{ln}\mathrm{\Delta })\gamma ^m\eta _1`$ (69)
$`0`$ $`=\beta ^{}\mathrm{\Delta }^1W^{}\eta _{1+}+\alpha {\displaystyle \frac{me^{5\varphi /4}}{16}}\eta _{1+}+5i\alpha f{\displaystyle \frac{e^{\varphi /4}}{32}}\eta _{1+}+\alpha {\displaystyle \frac{me^{3\varphi /4}}{32}}B_{mn}^{}\gamma ^{mn}\eta _{1+}`$
$`+\beta {\displaystyle \frac{e^{\varphi /2}}{96}}H_{mnp}\gamma ^{mnp}\eta _1\alpha {\displaystyle \frac{e^{\varphi /4}}{256}}G_{mnpq}\gamma ^{mnpq}\eta _{1+}+{\displaystyle \frac{1}{2}}\beta _m(\mathrm{ln}\mathrm{\Delta })\gamma ^m\eta _1,`$ (70)
from the noncompact piece. Note that these equations are complex. Similarly from the dilatino we obtain
$`0`$ $`={\displaystyle \frac{1}{2}}\alpha ^{}_m\varphi \gamma ^m\eta _1\alpha ^{}{\displaystyle \frac{e^{\varphi /2}}{24}}H_{mnp}\gamma ^{mnp}\eta _1\beta ^{}{\displaystyle \frac{5me^{5\varphi /4}}{4}}\eta _{1+}`$
$`+i\beta ^{}f{\displaystyle \frac{e^{\varphi /4}}{8}}\eta _{1+}\beta ^{}{\displaystyle \frac{3me^{3\varphi /4}}{8}}B_{mn}^{}\gamma ^{mn}\eta _{1+}\beta ^{}{\displaystyle \frac{e^{\varphi /4}}{192}}G_{mnpq}\gamma ^{mnpq}\eta _{1+}`$ (71)
$`0`$ $`={\displaystyle \frac{1}{2}}\beta _m\varphi \gamma ^m\eta _1+\beta {\displaystyle \frac{e^{\varphi /2}}{24}}H_{mnp}\gamma ^{mnp}\eta _1+\alpha {\displaystyle \frac{5me^{5\varphi /4}}{4}}\eta _{1+}`$
$`+i\alpha f{\displaystyle \frac{e^{\varphi /4}}{8}}\eta _{1+}\alpha {\displaystyle \frac{3me^{3\varphi /4}}{8}}B_{mn}^{}\gamma ^{mn}\eta _{1+}+\alpha {\displaystyle \frac{e^{\varphi /4}}{192}}G_{mnpq}\gamma ^{mnpq}\eta _{1+}.`$ (72)
A second set of conditions follows from the second supersymmetry (65). These can be obtained from the ones above by substituting $`(\alpha ,\beta ,\eta _1)(\gamma ,\delta ,\eta _2)`$.
### 6.2 Analysis of the conditions
Let us consider the scalar component of the supersymmetry equations first. It is straightforward to show that if there exists a point $`y_0`$ in $`X_6`$ such that $`|\alpha (y_0)||\beta (y_0)|`$ or $`|\gamma (y_0)||\delta (y_0)|`$, equations (69-72) (and the ones obtained from them by substituting $`(\alpha ,\beta ,\eta _1)(\gamma ,\delta ,\eta _2)`$) imply that $`m,f,W=0`$. I.e. the space $`M_{1,3}`$ reduces to Minkowski, which is contrary to our assumption. Hence, up to phases which can be absorbed in the normalizations of $`\eta _{1,2}`$ we can take:
$`\alpha `$ $`=\beta `$
$`\gamma `$ $`=\delta ,`$ (73)
at each point in $`X_6`$. Taking (73) into account, it is useful to note that the supersymmetry conditions (67, 68), as well as the ones obtained from them by substituting $`(\alpha ,\eta _1)(\gamma ,\eta _2)`$, are equivalent to the following set of equations:
$`0`$ $`=_m\eta _{1+}+_m\mathrm{ln}|\alpha |\eta _{1+}+3if{\displaystyle \frac{e^{\varphi /4}}{32}}\gamma _m\eta _1+{\displaystyle \frac{me^{3\varphi /4}}{32}}B_{np}^{}(\gamma _m{}_{}{}^{np}14\delta _m{}_{}{}^{n}\gamma _{}^{p})\eta _1`$ (74)
$`0`$ $`=_m\mathrm{ln}\left({\displaystyle \frac{\alpha }{|\alpha |}}\right)\eta _{1+}+{\displaystyle \frac{e^{\varphi /2}}{96}}H_{npq}(\gamma _m{}_{}{}^{npq}9\delta _m{}_{}{}^{n}\gamma _{}^{pq})\eta _{1+}{\displaystyle \frac{me^{5\varphi /4}}{16}}\gamma _m\eta _1`$
$`+{\displaystyle \frac{e^{\varphi /4}}{256}}G_{npqr}(\gamma _m{}_{}{}^{npqr}{\displaystyle \frac{20}{3}}\delta _m{}_{}{}^{n}\gamma _{}^{pqr})\eta _1`$ (75)
and
$`0`$ $`=_m\eta _{2+}+_m\mathrm{ln}|\gamma |\eta _{2+}+3if{\displaystyle \frac{e^{\varphi /4}}{32}}\gamma _m\eta _2+{\displaystyle \frac{me^{3\varphi /4}}{32}}B_{np}^{}(\gamma _m{}_{}{}^{np}14\delta _m{}_{}{}^{n}\gamma _{}^{p})\eta _2`$ (76)
$`0`$ $`=_m\mathrm{ln}\left({\displaystyle \frac{\gamma }{|\gamma |}}\right)\eta _{2+}+{\displaystyle \frac{e^{\varphi /2}}{96}}H_{npq}(\gamma _m{}_{}{}^{npq}9\delta _m{}_{}{}^{n}\gamma _{}^{pq})\eta _{2+}{\displaystyle \frac{me^{5\varphi /4}}{16}}\gamma _m\eta _2`$
$`+{\displaystyle \frac{e^{\varphi /4}}{256}}G_{npqr}(\gamma _m{}_{}{}^{npqr}{\displaystyle \frac{20}{3}}\delta _m{}_{}{}^{n}\gamma _{}^{pqr})\eta _2.`$ (77)
Let us first analyze the supersymmetry equations (75, 77), (69-72) and the ones obtained from them by $`(\alpha ,\eta _1)(\gamma ,\eta _2)`$, considering each irreducible $`SU(2)`$ representation in turn. The decompositions of all antisymmetric tensors in terms of irreducible $`SU(2)`$ representations can be found in appendix C. One can show that the solution is equivalent to the following conditions:
The $`\mathrm{𝟏}`$
$`m`$ $`=0`$
$`W`$ $`={\displaystyle \frac{i\mathrm{\Delta }}{6}}\left({\displaystyle \frac{\alpha }{|\alpha |}}\right)^2fe^{\varphi /4}.`$ (78)
$`mb_1`$ $`={\displaystyle \frac{f}{6}}e^{\varphi /2}`$
$`mb_2`$ $`={\displaystyle \frac{4i}{3}}e^{3\varphi /4}d_K\varphi `$
$`mb_3`$ $`=0,`$ (79)
$`g_1=g_2=g_3=0,`$ (80)
$`h_1=h_2=h_3=0,`$ (81)
$`{\displaystyle \frac{\alpha }{|\alpha |}}`$ $`=\pm {\displaystyle \frac{\gamma }{|\gamma |}}`$
$`d_K\left({\displaystyle \frac{\alpha }{|\alpha |}}\right)=`$ $`d_K^{}\left({\displaystyle \frac{\alpha }{|\alpha |}}\right)=0,`$ (82)
$`d_K\mathrm{ln}\mathrm{\Delta }`$ $`={\displaystyle \frac{1}{12}}d_K\varphi `$
$`d_K\varphi `$ $`=d_K^{}\varphi .`$ (83)
The first line of (83) together with the second line of (78) and the fact that $`W`$ is constant, imply:
$`d_K\mathrm{ln}f={\displaystyle \frac{1}{6}}d_K\varphi .`$ (84)
The $`\mathrm{𝟐}`$
$`m\stackrel{~}{b}_{1m}`$ $`={\displaystyle \frac{4i}{3}}e^{3\varphi /4}\stackrel{~}{}_m^+\varphi `$
$`m\stackrel{~}{b}_{2m}`$ $`={\displaystyle \frac{4i}{3}}e^{3\varphi /4}\stackrel{~}{}_m^{}\varphi ,`$ (85)
$`\stackrel{~}{h}_{1m}=\stackrel{~}{h}_{2m}=0,`$ (86)
$`\stackrel{~}{g}_{1m}=\stackrel{~}{g}_{2m}=0,`$ (87)
$`\stackrel{~}{}_m\left({\displaystyle \frac{\alpha }{|\alpha |}}\right)`$ $`=0`$
$`\stackrel{~}{}_m\mathrm{ln}\mathrm{\Delta }`$ $`={\displaystyle \frac{1}{12}}\stackrel{~}{}_m\varphi .`$ (88)
The $`\mathrm{𝟑}`$
$`\stackrel{~}{h}_{mn}=\stackrel{~}{g}_{mn}=0.`$ (89)
The relations derived so far imply $`\mathrm{\Delta }=\mathrm{constant}\times e^{\varphi /12}`$, $`f=\mathrm{constant}\times e^{\varphi /6}`$, as well as $`H=0`$, $`G=fdVol_4`$, where $`dVol_4`$ is the volume element of $`M_{1,3}`$ in the warped metric. It then follows from the Bianchi identity (7) for the $`G`$ field that $`\varphi =\mathrm{constant}`$.
To summarize the conditions so far:
$`m`$ $`=0`$
$`W`$ $`={\displaystyle \frac{i\mathrm{\Delta }}{6}}\left({\displaystyle \frac{\alpha }{|\alpha |}}\right)^2fe^{\varphi /4}`$
$`{\displaystyle \frac{\gamma }{|\gamma |}}`$ $`=\pm {\displaystyle \frac{\alpha }{|\alpha |}}`$
$`{\displaystyle \frac{\alpha }{|\alpha |}},\mathrm{\Delta },\varphi ,f`$ $`=\mathrm{constant}.`$ (90)
In addition, in form notation,
$`F`$ $`=\stackrel{~}{f}{\displaystyle \frac{i}{6}}fe^{\varphi /2}KK^{}`$
$`H`$ $`=0`$
$`G`$ $`=fdVol_4.`$ (91)
Note that we have taken (8) and the fact that $`m=0`$ into account, and we have set $`m\stackrel{~}{b}_{mn}=\frac{1}{2}\stackrel{~}{f}_{mn}`$.
Next we turn to the equations (74,76). The fact that $`\eta _{1,2}`$ are unimodular implies $`(\eta _1^+\eta _1)=0`$ and $`(\eta _2^+\eta _2)=0`$ which, taking (74,76) into account, can be seen to be equivalent to $`|\alpha |,|\gamma |=\mathrm{constant}`$. Together with (90) this implies
$`\alpha ,\gamma =\mathrm{constant}.`$ (92)
In addition, the orthogonality of $`\eta _{1,2}`$ implies $`(\eta _1^+\eta _2)=0`$ which, taking (74,76) into account, leads to the condition
$`f=0.`$ (93)
Taking (90) into account, this implies $`W=0`$ and $`M_{1,3}`$ reduces to Minkowski space<sup>8</sup><sup>8</sup>8It is not difficult to see that in addition the equations of motion impose $`\stackrel{~}{f}=0`$ and therefore all fluxes are zero.. Of course, this is contrary to our assumption of a (warped) $`AdS_4`$ vacuum. We are therefore led to the conclusion that there are no $`𝒩=2`$ solutions of type IIA supergravity satisfying our requirements.
Acknowledgements: We are grateful to C. Jeschek for valuable discussions. This work is supported in part by the EU-RTN network Constituents, Fundamental Forces and Symmetries of the Universe (MRTN-CT-2004-005104).
## Appendix A Fierz identities
Using definitions (23,24) we find
$`\eta _1^\alpha \eta _{1+}^\beta `$ $`={\displaystyle \frac{1}{4}}(P_{}C^1)^{\alpha \beta }+{\displaystyle \frac{i}{8}}J_{mn}(P_{}\gamma ^{mn}C^1)^{\alpha \beta }`$
$`\eta _{1+}^\alpha \eta _{1+}^\beta `$ $`={\displaystyle \frac{1}{48}}\mathrm{\Omega }_{mnp}(P_+\gamma ^{mnp}C^1)^{\alpha \beta }`$
$`\eta _1^\alpha \eta _1^\beta `$ $`={\displaystyle \frac{1}{48}}\mathrm{\Omega }_{mnp}^{}(P_{}\gamma ^{mnp}C^1)^{\alpha \beta }`$ (94)
and similarly for $`\eta _2\eta _2`$, by replacing $`(J,\mathrm{\Omega })(J^{},\mathrm{\Omega }^{})`$. Moreover for $`\eta _1\eta _2`$ we have
$`\eta _1^\alpha \eta _{2+}^\beta `$ $`={\displaystyle \frac{i}{8}}\omega _{mn}^{}(P_{}\gamma ^{mn}C^1)^{\alpha \beta }`$
$`\eta _2^\alpha \eta _{1+}^\beta `$ $`={\displaystyle \frac{i}{8}}\omega _{mn}(P_{}\gamma ^{mn}C^1)^{\alpha \beta }`$
$`\eta _{1+}^\alpha \eta _{2+}^\beta `$ $`={\displaystyle \frac{1}{4}}K_m(P_+\gamma ^mC^1)^{\alpha \beta }{\displaystyle \frac{1}{48}}\stackrel{~}{\mathrm{\Omega }}_{mnp}(P_+\gamma ^{mnp}C^1)^{\alpha \beta }`$
$`\eta _1^\alpha \eta _2^\beta `$ $`={\displaystyle \frac{1}{4}}K_m^{}(P_{}\gamma ^mC^1)^{\alpha \beta }+{\displaystyle \frac{1}{48}}\stackrel{~}{\mathrm{\Omega }}_{mnp}^{}(P_{}\gamma ^{mnp}C^1)^{\alpha \beta },`$ (95)
where
$`\stackrel{~}{\mathrm{\Omega }}_{mnp}`$ $`:=(\eta _2^+\gamma _{mnp}\eta _{1+}).`$ (96)
The latter is imaginary self-dual
$`\stackrel{~}{\mathrm{\Omega }}_{mnp}={\displaystyle \frac{i}{6}}\sqrt{\rho _6}\epsilon _{mnpqrs}\stackrel{~}{\mathrm{\Omega }}^{qrs}`$ (97)
and obeys
$`\stackrel{~}{\mathrm{\Omega }}\omega =\stackrel{~}{\mathrm{\Omega }}\omega ^{}=0.`$ (98)
As follows from (95), the two globally defined spinors are related via
$`\eta _{2+}={\displaystyle \frac{1}{2}}K^m\gamma _m\eta _1.`$ (99)
We also note the following relations,
$`0`$ $`=(\mathrm{\Pi }^+)_m{}_{}{}^{n}\gamma _{n}^{}\eta _1`$
$`\gamma _{mn}\eta _{1+}`$ $`=iJ_{mn}\eta _{1+}+{\displaystyle \frac{1}{2}}\mathrm{\Omega }_{mnp}\gamma ^p\eta _1`$
$`\gamma _{mnp}\eta _1`$ $`=3iJ_{[mn}\gamma _{p]}\eta _1\mathrm{\Omega }_{mnp}^{}\eta _{1+}.`$ (100)
A useful formula following from (99, 100) is
$`\gamma ^m\eta _2=K^m\eta _{1+}+{\displaystyle \frac{i}{2}}\omega _{mn}\gamma ^n\eta _1.`$ (101)
## Appendix B $`SU(2)`$ structure
Here we give some further useful relations pertaining to the $`SU(2)`$ structure.
It follows from (100) that
$`0`$ $`=(\stackrel{~}{\mathrm{\Pi }}^+)_m{}_{}{}^{n}\gamma _{n}^{}\eta _1,`$ (102)
where
$`(\stackrel{~}{\mathrm{\Pi }}^\pm )_{mk}:={\displaystyle \frac{1}{2}}(\stackrel{~}{\rho }_{mk}i\stackrel{~}{J}_{mk})`$ (103)
and
$`\stackrel{~}{\rho }_{mk}:=\rho _{mk}{\displaystyle \frac{1}{2}}(K_mK_k^{}+K_m^{}K_k).`$ (104)
Note that
$`K^m\stackrel{~}{\rho }_{mk}=0.`$ (105)
Some further useful identities are
$`\stackrel{~}{J}_{mn}\stackrel{~}{J}^n_k`$ $`=\stackrel{~}{\rho }_{mk}`$
$`\stackrel{~}{J}_m{}_{}{}^{n}\omega _{nk}^{}`$ $`=i\omega _{mk}`$
$`\omega _{mn}\omega ^{nk}`$ $`=4(\stackrel{~}{\mathrm{\Pi }}^+)_m^k`$
$`\omega _{mn}\omega ^{ij}`$ $`=8(\stackrel{~}{\mathrm{\Pi }}^+)_{[m}{}_{}{}^{i}(\stackrel{~}{\mathrm{\Pi }}^+)_{n]}^{}^j`$
$`(\stackrel{~}{\mathrm{\Pi }}^+)_m^k`$ $`=(\mathrm{\Pi }^+)_m{}_{}{}^{k}{\displaystyle \frac{1}{2}}K_mK^k.`$ (106)
## Appendix C $`SU(2)`$ tensor decompositions
In terms of the $`SU(2)`$ structure, the form fields of IIA supergravity decompose as follows.
Two-form
$`B_{mn}^{}=b_{mn}+b_{[m}K_{n]}+b_{[m}^{}K_{n]}^{}+ib_1K_{[m}K_{n]}^{},`$ (107)
where
$`K^ib_{im}=K^ib_i=K^ib_i=0`$ (108)
and
$`K^iB_{im}^{}`$ $`=b_m^{}ib_1K_m`$
$`K^iK^jB_{ij}^{}`$ $`=2ib_1.`$ (109)
Note that $`b_1`$ is real. We can further decompose
$`b_{mn}=\stackrel{~}{b}_{mn}+{\displaystyle \frac{1}{8}}\omega _{mn}^{}b_2+{\displaystyle \frac{1}{8}}\omega _{mn}b_2^{}+{\displaystyle \frac{1}{4}}\stackrel{~}{J}_{mn}b_3,`$ (110)
where $`\stackrel{~}{b}_{mn}`$ is $`(1,1)`$ and traceless with respect to $`\stackrel{~}{J}_{mn}`$, i.e. it transforms in the $`\mathrm{𝟑}`$ of $`SU(2)`$. The scalar $`b_2`$ is complex whereas $`b_3`$ is real. We have
$`b_2`$ $`=\omega ^{mn}b_{mn}`$
$`b_3`$ $`=\stackrel{~}{J}^{mn}b_{mn}.`$ (111)
Finally,
$`b_m={\displaystyle \frac{1}{4}}\omega _m^{}{}_{}{}^{i}\stackrel{~}{b}_{1i}^{}{\displaystyle \frac{1}{4}}\omega _m{}_{}{}^{i}\stackrel{~}{b}_{2i}^{},`$ (112)
where $`(\mathrm{\Pi }^{})_m{}_{}{}^{n}\stackrel{~}{b}_{1n}^{}=(\mathrm{\Pi }^+)_m{}_{}{}^{n}\stackrel{~}{b}_{2n}^{}=0`$. Both $`\stackrel{~}{b}_{1i}`$, $`\stackrel{~}{b}_{2i}`$ transform in the $`\mathrm{𝟐}`$ of $`SU(2)`$. We have
$`\stackrel{~}{b}_{1i}`$ $`=\omega _m{}_{}{}^{n}b_{n}^{}`$
$`\stackrel{~}{b}_{2i}`$ $`=\omega _m^{}{}_{}{}^{n}b_{n}^{}.`$ (113)
Three-form
$`H_{mnp}=h_{mnp}+h_{[mn}K_{p]}+h_{[mn}^{}K_{p]}^{}+ih_{[m}K_nK_{p]}^{},`$ (114)
where
$`K^ih_{imn}=K^ih_{im}=K^ih_{im}=K^ih_i=0`$ (115)
and
$`K^iH_{imn}`$ $`={\displaystyle \frac{2}{3}}h_{mn}^{}+{\displaystyle \frac{2i}{3}}h_{[m}K_{n]}`$
$`K^iK^jH_{ijm}`$ $`={\displaystyle \frac{2i}{3}}h_m.`$ (116)
Note that $`h_m`$ is real whereas $`h_{mn}`$ is complex. We can further decompose
$`h_{mnp}={\displaystyle \frac{3}{32}}\omega _{[mn}\omega _{p]}^{}{}_{}{}^{i}\stackrel{~}{h}_{1i}^{}{\displaystyle \frac{3}{32}}\omega _{[mn}^{}\omega _{p]}{}_{}{}^{i}\stackrel{~}{h}_{1i}^{},`$ (117)
where $`(\mathrm{\Pi }^{})_m{}_{}{}^{n}\stackrel{~}{h}_{1n}^{}=0`$. We have
$`\stackrel{~}{h}_{1m}`$ $`=\omega _m{}_{}{}^{i}\omega _{}^{jk}h_{ijk}.`$ (118)
Moreover
$`h_{mn}=\stackrel{~}{h}_{mn}+{\displaystyle \frac{1}{8}}\omega _{mn}^{}h_1+{\displaystyle \frac{1}{8}}\omega _{mn}h_2+{\displaystyle \frac{1}{4}}\stackrel{~}{J}_{mn}h_3,`$ (119)
where $`\stackrel{~}{h}_{mn}`$ is complex and it is $`(1,1)`$ and traceless with respect to $`\stackrel{~}{J}_{mn}`$. The scalars $`h_{1,2,3}`$ are complex. We have
$`h_1`$ $`=\omega ^{mn}h_{mn}`$
$`h_2`$ $`=\omega ^{mn}h_{mn}`$
$`h_3`$ $`=\stackrel{~}{J}^{mn}h_{mn}.`$ (120)
Finally,
$`h_m={\displaystyle \frac{1}{4}}\omega _m^{}{}_{}{}^{i}\stackrel{~}{h}_{2i}^{}{\displaystyle \frac{1}{4}}\omega _m{}_{}{}^{i}\stackrel{~}{h}_{2i}^{},`$ (121)
where $`(\mathrm{\Pi }^{})_m{}_{}{}^{n}\stackrel{~}{h}_{2n}^{}=0`$. We have
$`\stackrel{~}{h}_{2i}`$ $`=\omega _m{}_{}{}^{n}h_{n}^{}.`$ (122)
Four-form
$`G_{mnpq}=g_{mnpq}+g_{[mnp}K_{q]}+g_{[mnp}^{}K_{q]}^{}+ig_{[mn}K_pK_{q]}^{},`$ (123)
where
$`K^ig_{imnp}=K^ig_{imn}=K^ig_{imn}=K^ig_{im}=0`$ (124)
and
$`K^iG_{imnp}`$ $`={\displaystyle \frac{1}{2}}g_{mnp}^{}{\displaystyle \frac{i}{2}}g_{[mn}K_{p]}`$
$`K^iK^jG_{ijmn}`$ $`={\displaystyle \frac{i}{3}}g_{mn}.`$ (125)
Note that $`g_{mnpq}`$, $`g_{mn}`$ are real whereas $`g_{mnp}`$ is complex. We can further decompose
$`g_{mnpq}={\displaystyle \frac{3}{8}}\stackrel{~}{J}_{[mn}\stackrel{~}{J}_{pq]}g_1,`$ (126)
where the scalar $`g_1`$ is real. We have
$`g_1=\stackrel{~}{J}^{mn}\stackrel{~}{J}^{pq}g_{mnpq}.`$ (127)
Moreover
$`g_{mnp}={\displaystyle \frac{3}{32}}\omega _{[mn}\omega _{p]}^{}{}_{}{}^{i}\stackrel{~}{g}_{1i}^{}{\displaystyle \frac{3}{32}}\omega _{[mn}^{}\omega _{p]}{}_{}{}^{i}\stackrel{~}{g}_{2i}^{},`$ (128)
where $`(\mathrm{\Pi }^{})_m{}_{}{}^{n}\stackrel{~}{g}_{1n}^{}=(\mathrm{\Pi }^+)_m{}_{}{}^{n}\stackrel{~}{g}_{2n}^{}=0`$. We have
$`\stackrel{~}{g}_{1m}`$ $`=\omega _m{}_{}{}^{i}\omega _{}^{jk}g_{ijk}`$
$`\stackrel{~}{g}_{2m}`$ $`=\omega _m^{}{}_{}{}^{i}\omega _{}^{jk}g_{ijk}.`$ (129)
Finally,
$`g_{mn}=\stackrel{~}{g}_{mn}+{\displaystyle \frac{1}{8}}\omega _{mn}^{}g_2+{\displaystyle \frac{1}{8}}\omega _{mn}g_2^{}+{\displaystyle \frac{1}{4}}\stackrel{~}{J}_{mn}g_3,`$ (130)
where $`\stackrel{~}{g}_{mn}`$ is real and it is traceless with respect to $`\stackrel{~}{J}_{mn}`$. The scalar $`g_2`$ is complex whereas $`g_3`$ is real. We have
$`g_2`$ $`=\omega ^{mn}g_{mn}`$
$`g_3`$ $`=\stackrel{~}{J}^{mn}g_{mn}.`$ (131)
## Appendix D $`SU(2)`$ supersymmetry reduction
Using (99) and the decompositions of section C, it follows that conditions (67, 68) can be cast in the form
$`U_m\eta _{1+}+U_{mn}\gamma ^n\eta _1=0,`$ (132)
whereas conditions (69-72) can be written as
$`V\eta _{1+}+V_m\gamma ^m\eta _1=0,`$ (133)
for some $`U_m`$, $`U_{mn}`$, $`V`$, $`V_m`$. The explicit expressions for the $`U`$’s and $`V`$’s can be readily read off from the following decompositions in terms of irreducible $`SU(2)`$ representations:
Two-form
$`(\gamma _m{}_{}{}^{np}B_{np}^{}`$ $`14\gamma ^pB_{mp}^{})\eta _1=\{iK_m^{}b_2^{}2i\stackrel{~}{b}_{2m}\}\eta _{1+}`$
$`+\{(2b_13b_3)\stackrel{~}{J}_{mn}14\stackrel{~}{b}_{mn}2i\stackrel{~}{J}_m{}_{}{}^{i}\stackrel{~}{b}_{in}^{}{\displaystyle \frac{3}{2}}b_2^{}\omega _{mn}`$
$`{\displaystyle \frac{3}{2}}K_m\omega _n{}_{}{}^{i}\stackrel{~}{b}_{2i}^{}2K_m^{}\omega _n{}_{}{}^{i}\stackrel{~}{b}_{1i}^{}+K_n(2\omega _m^{}{}_{}{}^{i}\stackrel{~}{b}_{1i}^{}+{\displaystyle \frac{3}{2}}\omega _m{}_{}{}^{i}\stackrel{~}{b}_{2i}^{})+iK_nK_m^{}(7b_1{\displaystyle \frac{1}{2}}b_3)\}\gamma ^n\eta _1,`$ (134)
$`(\gamma _m{}_{}{}^{np}B_{np}^{}`$ $`14\gamma ^pB_{mp}^{})\eta _2=\{iK_m^{}(14b_1+b_3)+3\omega _m^{}{}_{}{}^{i}\stackrel{~}{b}_{1i}^{}+4\omega _m{}_{}{}^{i}\stackrel{~}{b}_{2i}^{}\}\eta _{1+}`$
$`+\{3b_2\stackrel{~}{J}_{mn}+(b_1+{\displaystyle \frac{3}{2}}b_3)\omega _{mn}i\omega _m{}_{}{}^{i}\stackrel{~}{b}_{in}^{}7i\omega _n{}_{}{}^{i}\stackrel{~}{b}_{im}^{}`$
$`3iK_m\stackrel{~}{b}_{1n}4iK_m^{}\stackrel{~}{b}_{2n}^{}iK_n\stackrel{~}{b}_{1m}{\displaystyle \frac{i}{2}}b_2K_nK_m^{}\}\gamma ^n\eta _1`$ (135)
and
$`B_{mn}^{}\gamma ^{mn}\eta _{1+}=i(b_3+2b_1)\eta _{1+}+\left\{i\stackrel{~}{b}_{2m}^{}{\displaystyle \frac{i}{2}}b_2K_m\right\}\gamma ^m\eta _1,`$ (136)
$`B_{mn}^{}\gamma ^{mn}\eta _{2+}=ib_2^{}\eta _{1+}+\left\{{\displaystyle \frac{1}{2}}\omega _m{}_{}{}^{i}\stackrel{~}{b}_{1i}^{}+iK_m({\displaystyle \frac{1}{2}}b_3b_1)\right\}\gamma ^m\eta _1.`$ (137)
Three-form
$`H_{npq}(\gamma _m{}_{}{}^{npq}`$ $`9\delta _m{}_{}{}^{n}\gamma _{}^{pq})\eta _{1+}=\{3\stackrel{~}{h}_{1m}+{\displaystyle \frac{3}{2}}\stackrel{~}{h}_{1m}^{}+2i\omega _m{}_{}{}^{i}\stackrel{~}{h}_{2i}^{}+i\omega _m^{}{}_{}{}^{i}\stackrel{~}{h}_{2i}^{}4ih_3K_m2ih_3^{}K_m^{}\}\eta _{1+}`$
$`+\{3K_m\stackrel{~}{h}_{2n}+K_n\stackrel{~}{h}_{2m}{\displaystyle \frac{3i}{8}}K_m\omega _n{}_{}{}^{i}\stackrel{~}{h}_{1i}^{}{\displaystyle \frac{9i}{8}}K_n\omega _m{}_{}{}^{i}\stackrel{~}{h}_{1i}^{}+2ih_1K_mK_n+ih_2^{}K_nK_m^{}`$
$`2h_2^{}\stackrel{~}{J}_{mn}+h_3^{}\omega _{mn}2i\omega _m{}_{}{}^{j}\stackrel{~}{h}_{jn}^{}6i\omega _n{}_{}{}^{j}\stackrel{~}{h}_{jm}^{}\}\gamma ^n\eta _1,`$ (138)
$`H_{npq}(\gamma _m{}_{}{}^{npq}`$ $`9\delta _m{}_{}{}^{n}\gamma _{}^{pq})\eta _{2+}=\{2\stackrel{~}{h}_{2m}^{}+{\displaystyle \frac{9i}{4}}\omega _m^{}{}_{}{}^{i}\stackrel{~}{h}_{1i}^{}4iK_mh_22iK_m^{}h_1^{}\}\eta _{1+}`$
$`+\{{\displaystyle \frac{3}{2}}K_n\stackrel{~}{h}_{1m}^{}{\displaystyle \frac{3}{4}}K_n\stackrel{~}{h}_{1m}{\displaystyle \frac{3}{4}}K_m\stackrel{~}{h}_{1n}+{\displaystyle \frac{3i}{2}}K_m\omega _n{}_{}{}^{i}\stackrel{~}{h}_{2i}^{}{\displaystyle \frac{i}{2}}K_n\omega _m{}_{}{}^{i}\stackrel{~}{h}_{2i}^{}iK_n\omega _m^{}{}_{}{}^{i}\stackrel{~}{h}_{2i}^{}`$
$`2ih_3K_mK_nih_3^{}K_nK_m^{}+2h_3^{}\stackrel{~}{J}_{mn}+4i\stackrel{~}{J}_m{}_{}{}^{j}\stackrel{~}{h}_{jn}^{}+\omega _{mn}h_1^{}+12\stackrel{~}{h}_{mn}^{}\}\gamma ^n\eta _1`$ (139)
and
$`H_{mnp}\gamma ^{mnp}\eta _1=2ih_2\eta _{1+}+\left\{{\displaystyle \frac{i}{2}}\omega _m{}_{}{}^{i}\stackrel{~}{h}_{2i}^{}{\displaystyle \frac{3}{4}}\stackrel{~}{h}_{1m}ih_3K_m\right\}\gamma ^m\eta _1,`$ (140)
$`H_{mnp}\gamma ^{mnp}\eta _2=2ih_3\eta _{1+}+\left\{{\displaystyle \frac{3i}{8}}\omega _m{}_{}{}^{i}\stackrel{~}{h}_{1i}^{}\stackrel{~}{h}_{2m}ih_1K_m\right\}\gamma ^m\eta _1.`$ (141)
Four-form
$`(\gamma _m{}_{}{}^{npqr}G_{npqr}^{}`$ $`{\displaystyle \frac{20}{3}}\gamma ^{pqr}G_{mpqr})\eta _1=\{{\displaystyle \frac{10}{3}}K_m^{}g_2^{}{\displaystyle \frac{5i}{2}}\omega _m^{}{}_{}{}^{i}\stackrel{~}{g}_{1i}^{}\}\eta _{1+}`$
$`+\{i(5g_1+{\displaystyle \frac{2}{3}}g_3)\stackrel{~}{J}_{mn}+{\displaystyle \frac{20i}{3}}\stackrel{~}{g}_{mn}4\stackrel{~}{J}_m{}_{}{}^{i}\stackrel{~}{g}_{in}^{}+{\displaystyle \frac{i}{3}}g_2^{}\omega _{mn}`$
$`{\displaystyle \frac{1}{2}}K_m\stackrel{~}{g}_{1n}2K_m^{}\stackrel{~}{g}_{2n}^{}+K_n({\displaystyle \frac{1}{2}}\stackrel{~}{g}_{1m}2\stackrel{~}{g}_{2m})+K_nK_m^{}({\displaystyle \frac{5}{3}}g_3{\displaystyle \frac{3}{2}}g_1)\}\gamma ^n\eta _1,`$ (142)
$`(\gamma _m{}_{}{}^{npqr}G_{npqr}^{}`$ $`{\displaystyle \frac{20}{3}}\gamma ^{pqr}G_{mpqr})\eta _2=\{K_m^{}(3g_1+{\displaystyle \frac{10}{3}}g_3)4\stackrel{~}{g}_{1m}+\stackrel{~}{g}_{2m}\}\eta _{1+}`$
$`+\{{\displaystyle \frac{2i}{3}}g_2\stackrel{~}{J}_{mn}+({\displaystyle \frac{5i}{2}}g_1{\displaystyle \frac{i}{3}}g_3)\omega _{mn}2\omega _m{}_{}{}^{i}\stackrel{~}{g}_{in}^{}{\displaystyle \frac{10}{3}}\omega _n{}_{}{}^{i}\stackrel{~}{g}_{im}^{}`$
$`+{\displaystyle \frac{i}{4}}K_m\omega _n{}_{}{}^{i}\stackrel{~}{g}_{2i}^{}+iK_m^{}\omega _n{}_{}{}^{i}\stackrel{~}{g}_{1i}^{}{\displaystyle \frac{5i}{4}}K_n\omega _m{}_{}{}^{i}\stackrel{~}{g}_{2i}^{}+{\displaystyle \frac{5}{3}}g_2K_nK_m^{}\}\gamma ^n\eta _1`$ (143)
and
$`G_{mnpq}\gamma ^{mnpq}\eta _{1+}=(3g_1+2g_3)\eta _{1+}+\left\{{\displaystyle \frac{3i}{4}}\omega _m{}_{}{}^{i}\stackrel{~}{g}_{1i}^{}+g_2K_m\right\}\gamma ^m\eta _1,`$ (144)
$`G_{mnpq}\gamma ^{mnpq}\eta _{2+}=2g_2^{}\eta _{1+}+\left\{{\displaystyle \frac{3}{2}}\stackrel{~}{g}_{2m}^{}+K_m({\displaystyle \frac{3}{2}}g_1g_3)\right\}\gamma ^m\eta _1.`$ (145)
Finally, a derivative ($`_m`$) on $`X_6`$ will be decomposed as
$`_m=\stackrel{~}{}_m+{\displaystyle \frac{1}{2}}K_mK^n_n+{\displaystyle \frac{1}{2}}K_m^{}K^n_n,`$ (146)
so that
$`\iota _K\stackrel{~}{}=0.`$ (147)
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# An Overview of Complex Adaptive Systems
## 1 Basics of complex adaptive systems (CAS)
Definition (1): A complex adaptive system consists of inhomogeneous, interacting adaptive agents. Adaptive means capable of learning.
Definition (2): An emergent property of a CAS is a property of the system as a whole which does not exist at the individual elements (agents) level.
Typical examples are the brain, the immune system, the economy, social systems, ecology, insects swarm, etc..
Therefore to understand a complex system one has to study the system as a whole and not to decompose it into its constituents. This totalistic approach is against the standard reductionist one, which tries to decompose any system to its constituents and hopes that by understanding the elements one can understand the whole system.
## 2 Why should we study complex adaptive systems?
Most of living systems are CAS. Moreover they have intrinsic unpredictability which causes some ”seemingly wise” decisions to have harmful side effects. Therefore we should try to understand CAS to try to minimize such side effects. Here we give two examples of these side effects.
Mathematical models have played important roles in understanding the impact of vaccination programs. The complications of infectious diseases spread make the problem of predicting the impact of vaccinations a nonlinear problem. Sometimes a counter-intuitive result appears e.g. the threshold phenomena \[Edelstein-Keshet 1988\]. Here another example will be mentioned.
Several vaccination programs are known e.g. mass vaccination where all population is vaccinated, target vaccination where only a certain group is vaccinated.
If one tries to understand the expected impact of a vaccination program one should take the following points into account:
1. Vaccination is not perfect hence a probability of vaccination failure should be assumed.
2. Sometimes vaccination takes time to be effective.
3. Immunity is waning i.e. may be lost with time.
4. Long range contacts can play a significant role e.g. SARS (severe acute respiratory syndrome) has been transmitted between countries via air travellers.
Rubella is a mild viral infectious disease. Typically it is most dangerous when infecting a pregnant female where it has severe effects on the fetus. Once one gets it he (she) gets a life long immunity. There are several vaccination strategies for rubella \[Vynnycky et al 2003\]. The US policy is to vaccinate all two years old children. The UK policy is to vaccinate only 14-years old girls. Another strategy which is adopted in some underdeveloped countries is not to vaccinate at all. It has been found \[Jazbec et al 2003\] that in most cases the UK strategy is equal or better than the US one despite being cheaper.
An interesting situation arose when some countries adopted a private sector vaccination to MMR (Measles, Mumps and Rubella) \[Vynnycky et al 2003\]. It was expected that the number of Congenital Rubella Syndrome (CRS) will decrease. However it did not and in some countries (e.g. Greece and Costs Rica) it increased. The reason can be understood as follows: This vaccination to part of the population decreases the probability of contracting the disease at young age. Hence the number of susceptible individuals at adulthood increases. Consequently the probability of contracting the disease at adulthood increases. This is an example of the counterintuitive effects of some vaccination programs.
Another example for bad side effects is Lake Victoria \[Chu et al 2003\] where a new species called Nile perch was introduced expecting that it is more economically profitable. Yet the following results have appeared:
1. The local fishermen’s tools were not suitable for the new fish hence only large corporations benefited.
2. Due to its higher price the locals were unavailable to buy the new type.
3. The original fish used to eat the larva of mosquitoes but now mosquitoes’ numbers have increased significantly thus the quality of life of the locals have deteriorated!!
There are at least two sources for unpredictability in CAS. The first is the nonlinear interactions between its agents \[West 1990\]. The second is that CAS are open systems hence perturbation to one system may affect another related one e.g. perturbation to Lake Victoria affected the number of mosquitoes.
## 3 How to model a CAS?
The standard approaches are
1. Ordinary differential equations (ODE), difference equations and partial differential equations (PDE).
2. Cellular automata (CA) \[Ilachinski 2001\].
3. Evolutionary game theory \[Hofbauer and Sigmund 1998\].
4. Agent based models.
5. Networks \[Watts and Strogatz 1998\] etc..
6. Fractional calculus \[Stanislavsky 2000\].
Some of these approaches are included in \[Boccara 2004\].
The ODE and PDE approaches have some difficulties as follows \[Louzon et al 2003\]:
1. ODE and PDE assumes that local fluctuations have been smoothed out.
2. Typically they neglect correlations between movements of different species.
3. They assume instantaneous results of interactions.
Most biological systems show delay and do not satisfy the above assumptions. They concluded that a cellular automata (CA) \[Ilachinski 2001\] type system called microscopic simulation is more suitable to model complex biological systems. We agree that CA type systems are more suitable to model complex biological systems but such systems suffer from a main drawback namely the difficulty of obtaining analytical results. The known analytical results about CA type systems are very few compared to the known results about ODE and PDE. Some mathematical results about CA are given in the appendix.
Now we present a compromise i.e. a PDE which avoids the delay and the correlations drawbacks. It is called telegraph reaction diffusion equations \[Ahmed and Hassan 2000\]. To overcome the non-delay weakness in Fick’s law it is replaced by
$$J(x,t)+\tau \frac{J(x,t)}{t}=D\frac{c}{x},$$
(1)
where the flux $`J(x,t)`$ relaxes, with some given characteristic time constant $`\tau `$ and $`c`$ is the concentration of the diffusing substance. Combining Eq. (1) with the equation of continuity, one obtains the modified diffusion equation or the Telegraph equation:
$$\frac{c}{t}+\tau \frac{^2c}{x^2}=D\frac{^2c}{x^2}.$$
(2)
The corresponding Telegraph reaction diffusion (TRD) is given by
$$\tau \frac{^2c}{t^2}+\left(1\frac{\mathrm{d}f(c)}{\mathrm{d}c}\right)\frac{c}{t}=D\frac{^2c}{x^2}+f(c),$$
(3)
where $`f(c)`$ is a polynomial in $`c`$.
Another motivation for TRD comes from media with memory where the flux $`J`$ is related to the density $`c(x,t)`$ through a relaxation function $`K(t)`$ as follows
$$J(x,t)=_0^tK(t\stackrel{´}{t})\frac{c(x,\stackrel{´}{t})}{x}d\stackrel{´}{t}.$$
It can be shown \[Compte & Metzler 1997\] that, with a suitable choice for the kernel $`K(t)`$, the standard Telegraph equation is obtained.
A third motivation is that starting from discrete space time one does not obtain the standard diffusion equation but the telegraph equation \[Chopard and Droz 1991\].
Moreover it is known that TRD results from correlated random walk \[Diekmann et al, 2000\]. This supports the conclusion that Telegraph reaction diffusion equation is more suitable for modeling complex systems than the usual diffusion one.
## 4 The immune system as a complex system \[Segel and Cohen 2001, Ahmed and Hashish 2004\]
The emergent properties of the immune system (IS) included:
1. The ability to distinguish any substance (typically called antigen Ag) and determine whether it is damaging or not. If Ag is non-damaging (damaging) then, typically, IS tolerates it (responds to it).
2. If it decides to respond to it then IS determines whether to eradicate it or to contain it.
3. The ability to memorize most previously encountered Ag, which enables it to mount a more effective reaction in any future encounters. This is the basis of vaccination processes.
4. IS is complex thus it has a network structure.
5. The immune network is not homogeneous since there are effectors with many connections and others with low number of connections.
6. The Ag, which enters our bodies, has extremely wide diversity. Thus mechanisms have to exist to produce immune effectors with constantly changing random specificity to be able to recognize these Ag. Consequently IS is an adaptive complex system.
7. Having said that, one should notice that the wide diversity of IS contains the danger of autoimmunity (attacking the body). Thus mechanisms that limit autoimmunity should exist.
8. In addition to the primary clonal deletion mechanism, two further brilliant mechanisms exist: The first is that the IS network is a threshold or ”window” one i.e. no activation exists if the Ag quantity is too low or too high (This is called low and high zone tolerance).
9. Thus an auto reactive immune effector (i.e. an immune effector that attacks the body to which it belongs) will face so many self-antigens that it has to be suppressed due to the high zone tolerance mechanism.
10. Another mechanism against autoimmunity is the second signal given by antigen presenting cells (APC). If the immune effector is self reactive then, in most cases, it does not receive the second signal thus it becomes anergic.
11. Also long term memory can be explained by the phenomena of high and low zone tolerance where IS tolerates Ag if its quantity is too high or too low. So persisting Ag is possible and continuous activation of immune effectors may occur.
12. There is another possible explanation for long term memory using the immune system (Extremal Dynamics).
13. Thus design principles of IS can explain important phenomena of IS.
An interesting example is given by Matzinger \[Matzinger 2002\] where she argued that to prevent transplant rejection it may be more useful to design drugs that blocks signal II and not signal I (which the present drugs do). The reason is blocking signal II make the effectors (which originally were capable of recognizing the transplant) anergic while leaving the other immune effectors intact.
## 5 Conclusions
1. CAS should be studied as a whole hence reductionist point of view may not be reliable in some cases.
2. CAS are open with nonlinear local interactions hence:
1. Long range prediction is highly unlikely \[Strogatz 2000, Holmgren 1996\].
2. When studying a CAS take into consideration the effects of its perturbation on related systems e.g. perturbation of lake Victoria has affected mosquitoes’ numbers hence the locals quality of life. This is also relevant to the case of natural disasters where an earthquake at a city can cause a widespread power failure at other cities.
3. Expect side effects to any ”WISE” decision.
4. Mathematical and computer models may be helpful in reducing such side effects.
3. Optimization in CAS should be multi-objective and not single objective \[Collette and Siarry 2003\].
4. CAS are very difficult to control. Interference at highly connected sites may be a useful approach \[Dorogovtsev and Mendez 2004\]. The interlinked nature of CAS elements complicates both the unpredictability and controllability problems. It also plays an important role in innovations spread.
5. Memory effects should not be neglected in CAS. This lends more support for the proposed telegraph reaction diffusion Eq. (3). Also memory games have been studied \[Smale 1980, Ahmed and Hegazi 2000\]. Also delay and fractional calculus are relevant to CAS.
6. Mathematical topics motivated by CAS include ODE and PDE (non-autonomous, delayed, periodic coefficients, stability and persistence), multi-objective optimization (including biologically motivated methods e.g. Ant colony optimization, Extremal optimization, Genetic algorithm etc), difference equations, cellular automata, networks, fractional calculus, control (e.g. bounded delayed control of distributed systems), game theory, nonlinear dynamics and fuzzy mathematics.
Some of the mathematics motivated by CAS will be reviewed in the appendices.
## Acknowledgments
One of the authors (A. S. Hegazi) acknowledge the financial support of a Mansoura University grant.
## References
1. Ahmed E. and Hassan S.Z. (2000), ”On Diffusion in Some Biological and Economic Systems”, Z. Naturforsch. 55a, 669.
2. Ahmed E. and Hashish A.H. (2004), ”On Modeling the immune system as a complex system”, Theor. BioSci. (Accepted).
3. Ahmed E. and Hegazi A.S.”On Discrete Dynamical Systems Associated to Games”, Adv. Complex. Sys. 2, 423-429 (2000).
4. Ahmed E., Hegazi A.S. ”On Persistence and Stability of Some Spatially Inhomogeneous Systems”, J. Math. Analys. Appl. 268, 74-88 (2002).
5. Ahmed E., Hegazi A.S., Elgazzar A.S. and Yehia H.M., ”On Synchronization, Persistence and Seasonality in Some Inhomogeneous Models in Epidemics and Ecology”, Physica A 322, 155-168 (2003).
6. Ahmed E. Hegazi A.S. and Abdel-Hafiz A.T., ”On Multiobjective oligopoly” Nonlin. Dyn. Psych. Life Sci. 7, 203-217 (2003).
7. Ahmed, E. and Elgazzar, A.S., ”On Coordination and Continuous Hawk-Dove Games on Small-World Networks”, Eur. Phys. J. B 18, 159 (2000 a).
8. Ahmed, E. and Elgazzar, A.S., ”On the Dynamics of Local Hawk-Dove Game”, Int. J. Mod. Phys. C 11, 607 (2000 b).
9. Ahmed, E., Hegazi, A.S. and Elgazzar, A.S., ”An Epidemic Model on Small-World Networks and Ring Vaccination”, Int. J. Mod. Phys. C 13, 189 (2002).
10. Albert, R. and Barabási, A.L., ”Statistical mechanics of complex networks”, Rev. Mod. Phys. 74, 47 (2002).
11. Bagnoli F., Boccara N. and Rechtman R.(2002), ”Nature of phase transitions in probabilistic cellular automata with two absorbing states”, Lect. notes Computer sci. 2493, 249-258.
12. Barnett S. (1990),”Matrices”,Oxford Univ. Press, U.K.
13. Boccara N. (2004), ”Modeling complex systems”, Springer Publ., Berlin.
14. Chopard B. and Droz M. (1991), ”Cellular automata model for the diffusion equation”, J. Stat. Phys. 64, 859.
15. Chu D., Strand R. and Fjelland R. (2003), ”Theories of Complexity”, Complexity 8, 19.
16. Collette Y. and Siarry P. (2003), ”Multiobjective Optimization”, Springer, Berlin.
17. Compte A. and Metzle R.(1997), ”The generalized Cattaneo equation for the description of anomalous transport processes” J. Phys. A 30, 7277.
18. Diekmann O., R. Durrett, K. P. Hadeler, P. Maini, H. L. Smith, V. Capasso (eds.) (2000), ”Mathematics inspired by biology”, Springer Verlag, Berlin.
19. Dorogovtsev S.N. and Mendes J.F (2004), ”The shortest path to complex networks”, Cond-mat 0404593.
20. Domany E. and Kinzel W.(1984), Equivalence of CA and Ising models and directed percolation, Phys.Rev.Lett.53, 311.
21. Edelstein-Keshet L. (1988), ”Introduction to mathematical biology”, Random House, N.Y.
22. Elgazzar, A.S., ”Application of the Sznajd Sociophysics Model to Small-World Networks”, Int. J. Mod. Phys. C 12, 1537 (2001).
23. Elgazzar, A.S., ”A Model for the Evolution of Economic Systems in Social Networks”, Physica A 303, 543 (2002).
24. Erdös, P. and Rényi, A., ”On the Evolution of Random Graphs”, Publ. Math. Inst. Hung. Acad. Sci. 5, 17 (1960).
25. Holmgren R. (1996), ”A first course in discrete dynamical systems”, Springer, Berlin.
26. Hofbauer J. and Sigmund K. (1998), ”Evolutionary games and population dynamics”, Cambridge Univ. Press, U.K.
27. Ilachinski A.(2001), ”Cellular automata”, World Scientific Publ., Singapore.
28. Jazbec A., Delimar M. and Vrzic V.S. (2003), ”Simulation model of Rubella”, Appl. Math. Comp. (To Appear).
29. Jen E.(1990), ”Aperiodicity in 1-dimensional cellular automata”, Physica D 45, 3.
30. Kaneko K. (1993), ”Theory and application of coupled map lattices,” Wiley Pub.
31. Louzoun Y., Solomon S., Atlan H. and Cohen I.R. (2003), ”Proliferation and competition in discrete biological systems”, Bull. Math. Biol. 65, 375.
32. Martin O., Odlyzko A.M. and Wolfram S. (1984), ”Algebraic properties of cellular automata”, Comm. Math. Phys. 93, 219.
33. Matzinger P. (2002), ”The danger model”, Science 296, 12, pp. 301.
34. Milgram, S., ”The Small World Problem”, Psychol. Today 2, 60 (1967).
35. Rocco A. and West B.J. (1999), ”Fractional Calculus and the Evolution of Fractal Phenomena”, Physica A 265, 535.
36. Schonfisch B. and De Roos A.(1999), ”Synchronous and asynchronous updating in cellular automata”, Biosys. 51, 123-143.
37. Segel L.A. and Cohen I.R.(eds.)(2001), ”Design principles for the immune system and other distributed autonomous systems”, Oxford Univ. Press, U.K.
38. Smale S.(1980). ”The prisoners’ dilemma and dynamical systems associated to games”, Econometrica 48, 1617-1634.
39. Smith J.B.(2003), ”Complex systems”, CS 0303020.
40. Stanislavsky A.A. (2000), ”Memory effects and macroscopic manifestation of randomness”, Phys. Rev. E 61, 4752.
41. Stevens J.G., Rosenzweig R.E. and Cerkanowicz A.E.(1993), ”Transient and cyclic behavior of cellular automata with null boundary conditions”, J. Stat. Phys. 73, 159-174.
42. Strogatz S. (2001), ”Nonlinear dynamics and chaos”, Perseus Books Group
43. Tadaki S. (1994), ”Periodicity of cylindrical cellular automata, Nonlin. 9412003
44. Vynnycky E., Gay N.J. and Cutts F.T. (2003), ”The predicted impact of private sector MMR vaccination on the burden of Congenital Rubella Syndrome”, Vaccine 21, 2708.
45. Watts D.J. and Strogatz S.H. (1998), ”Collective dynamics of small world network”, Nature 393, 440.
46. West B.J. (1991), ”Fractal Physiology and chaos in medicine”, World Scientific Publ., Singapore.
## Appendix (1): Some mathematical results for one-dimensional cellular automata
Definition (3): A cellular automata consists of 4 components: A graph $`G`$, a set of states such that each site (vertex) of the graph has one of the possible states, a neighborhood set which assigns to each vertex a certain neighborhood and a transition function $`f`$ which defines the evolution of the state of each site as a function of the states of that site and those in its neighborhood.
We choose the set of possible states to be the ring $`Z(p)`$ i.e. the set of integers $`0,1,2,\mathrm{},p1`$, where addition is defined $`\mathrm{mod}p`$. The total number of sites is denoted by $`N`$. In most of the cases, we choose $`N,p`$ to be relatively prime. The set of states of the sites at a given time is called a configuration. We now restrict us to a one-dimensional space. Let $`x(j,t)`$ be the state of site $`j`$ at time $`t`$.
Definition (4): A finite initial configuration is one such that there are two natural numbers $`L,R`$ such that $`0<L<R<N`$, and $`x(j,0)=0`$ if $`j<L`$ or $`j>R`$.
Theorem (1) \[Jen 1990\]: If $`x(i,t),x(j,t),i<j`$ are two periodic sequences i.e. $`x(i,t)=x(i,t+p(i)),x(j,t)=x(j,t+p(j))`$, then for every $`k`$ such that $`i<k<j`$ then $`x(k,t)`$ is periodic.
Corollary (1) \[Jen 1990\]: If CA evolves according to the rule
$$x(i,t)=f(x(i1,t),x(i,t),x(i+1,t))\mathrm{mod}\mathrm{\hspace{0.33em}2},$$
(4)
such that $`0000,1000,0010`$, then for any finite initial configuration the system is temporarily periodic i.e. the sequence $`(x(i,t))`$ is periodic for all $`i`$ such that $`0<iN,0<T<t`$.
Proof. The fact that $`1000`$ implies that $`x(i,t)=0`$ for $`i>R`$, similarly $`x(i,t)=0`$ for $`i<L`$ for all $`t>0`$. Applying theorem (1) the result is proved.
In the case that $`f`$ in Eq. (4) is linear, one can use the methods of \[Stevens et al 1993, Tadaki 1994\] to get useful information about possible periodicity’s of the system. In this case the system can be written as
$$X(t+1)=UX(t),$$
(5)
where $`U`$ is called the evolution matrix. Then $`X(t)=U^tX(0)`$. In this case the asymptotic behavior of the system is governed by the characteristic polynomial of $`U`$ on the field $`Z(p)`$. Assuming periodic boundary conditions, the matrix $`U`$ is circulant matrix \[Barnett 1990\].
Let $`P(N,\lambda )`$ be the characteristic polynomial of the system (5) with $`N`$ sites, then typically it has the form
$$P(n,\lambda )=\lambda ^ad(n,\lambda ),d(n,0)=1.$$
(6)
If $`a>0`$, then the systems tends to a fixed configuration (which corresponds to a fixed point for discrete time continuous state dynamical systems). Reducing $`d(n,\lambda )`$ to its irreducible factors on the field of states then in most cases a cycle of length $`p^k1`$ exist for the system where $`k`$ is the degree of the irreducible factors.
As an example consider rule 90 \[Martin et al 1984\]
$$x(i,t+1)=x(i1,t)+x(i+1,t)\mathrm{mod}\mathrm{\hspace{0.33em}2}.$$
(7)
For $`N=5`$, we have $`P(5,\lambda )=\lambda (\lambda ^2+\lambda +1)^2\mathrm{mod}\mathrm{\hspace{0.33em}2}`$, hence the system may evolve to a fixed configuration (e,g, $`x(i,t)=0`$ for all $`t>T>0`$, for all $`0<iN`$). It can also evolve to a cycle of period $`3(=2^21)`$.
Similarly for $`N=9`$, $`P(9,\lambda )=\lambda (\lambda +1)^2(\lambda ^3+l+1)^2`$ on $`Z(2)`$. Hence this system may evolve to a fixed configuration or to a periodic one with period 7. For $`N=13`$, similar study implies that $`P(13,\lambda )=\lambda (\lambda ^6+\lambda ^5+\lambda ^4+\lambda +1)^2`$ on $`Z(2)`$. Hence fixed configurations and periodic ones with period 63 are expected. Such long periods may not be easy to find numerically. These results can be obtained using more elaborate methods \[Martin et al 1984\]; but the simplicity of the present approach is appealing.
Moreover it is directly applicable to nonlocal cases which have gained much attention after the pioneering work of Watts and Strogatz on small world network (SWN) \[Watts and Strogatz 1998\]. As an example consider the following system
$$x(i,t+1)=x(i1,t)+x(i+1,t)+x(i+k,t)\mathrm{mod}\mathrm{\hspace{0.33em}2},$$
(8)
where $`k`$ is fixed. Some of the characteristic polynomials P(N,$`\lambda `$,k) are:
$`P(11,0,\lambda )`$ $`=`$ $`\lambda ^{11}+\lambda ^{10}+\lambda ^5+\lambda ^4+\lambda +1,`$
$`P(11,3,\lambda )`$ $`=`$ $`\lambda ^{11}+\lambda ^9+\lambda ^7+\lambda ^6+\lambda ^5+\lambda ^4+\lambda +1,`$
$`P(11,1,\lambda )`$ $`=`$ $`\lambda ^{11}+\lambda ^8+\lambda ^7+\lambda ^5+\lambda ^2+1,`$ (9)
Hence we have the following proposition:
Proposition (1): a) The system (8) depends on $`k`$.
b) The asymptotic behavior of (8) contains the following: For $`N=11,k=3`$, no fixed configuration but a periodic one with period 1023.
Proof. a) For $`N=11,k=5`$, a homogeneous configuration is expected. This is not the case for $`N=11,k=0`$ or $`k=3`$.
b) Use the procedure explained before.
Typically updating of CA is synchronous. It is important to notice that other types of updating e.g. a uniform random asynchronous one (where only one site is chosen randomly and updated at each time step) gives other patterns \[Schonfisch and de Roos 1999\]. The following lemma is useful
Lemma (1): a) States which are stationary under synchronous updating are also stationary under asynchronous one.
b) If there is a site $`j`$ which is not updated for all time $`t>T>0`$ then stationary configuration with respect to asynchronous updating may not be so under synchronous one.
Proof. a) If $`f(x(1),x(2),\mathrm{},x(N))=(x(1),x(2),\mathrm{},x(N))`$ then $`f(j,x(j))=x(j)`$. This proves part a). Since site $`j`$ is not updated for $`t>T>0`$ then $`f(j,x(j))x(j)`$ can still belong to a homogeneous configuration for the asynchronous updating but not the homogeneous one. This proves b).
Loosely speaking patterns present in asynchronous updating are mostly present in synchronous one. Motivated by these results we study sequential CA e.g. the sequential rule 90 is
$$x(j,t+1)=x(j1,t+1)+x(j+1,t)\mathrm{mod}\mathrm{\hspace{0.33em}2}.$$
(10)
This can be written in the following equivalent form
$$x(j,t+1)=\underset{k=2}{\overset{j+2}{}}x(k,t)\mathrm{mod}\mathrm{\hspace{0.33em}2},$$
(11)
where free periodic boundary conditions are assumed. The characteristic polynomials of the system (7) are:
$`P(5,\lambda )`$ $`=`$ $`\lambda ^5+\lambda ^3,P(6,\lambda )=\lambda ^6+\lambda ^5+\lambda ^3,P(7,\lambda )=\lambda ^7,`$
$`P(13,\lambda )`$ $`=`$ $`\lambda ^7(\lambda ^3+\lambda ^2+1)^2.`$
Hence homogeneous configurations are expected for $`N=5,6,7`$. For $`N=13`$ a periodic configuration with period 7 is expected.
Studying the system (10) numerically showed that chaos (in the sense of sensitive dependence on initial conditions which is sometimes called damage spread) exists.
Proposition (2): Every initially finite configuration will evolve under the CA
$$x(j,t+1)=x(j1,t+1)x(j+1,t)\mathrm{mod}\mathrm{\hspace{0.33em}2},$$
(12)
into the zero configuration $`x(j,t)=0`$ for all $`j`$, $`0<jN`$, for all time $`t>T>0`$ where $`T<N`$.
Proof. We have
$$x(R,1)=x(R+1,0)x(R1,1).$$
But $`x(R+1,0)=0`$ by definition of initially finite configuration thus $`x(R,1)=0`$. Repeating for $`x(R1,2)`$, one gets $`x(R1,2)=0`$ and continue.
Now the above results are applied to two known examples. The first is Domany-Kinzel (DK) model \[Kinzel and Domany 1984\], which is given by:
$`\mathrm{If}x(j1,t)+x(j+1,t)`$ $`=`$ $`0\mathrm{then}x(j,t+1)=0.`$
$`\mathrm{If}x(j1,t)+x(j+1,t)`$ $`=`$ $`1\mathrm{then}x(j,t+1)=1,`$
$`\mathrm{with}\mathrm{probability}p_1.`$ (13)
$`\mathrm{If}x(j1,t)+x(j+1,t)`$ $`=`$ $`2\mathrm{then}x(j,t+1)=1,`$
$`\mathrm{with}\mathrm{probability}p_2.`$
where $`x(j,t)`$ are Boolean variables. For $`p_11`$, $`p_21`$, the system (13) corresponds to the CA
$$x(j,t+1)=x(j1,t)+x(j+1,t)+x(j1,t)x(j+1,t)\mathrm{mod}\mathrm{\hspace{0.33em}2}.$$
(14)
Proposition (3): Any finite initial configuration with two consecutive ones will tend to the homogeneous configuration $`x(j,t)=1`$ for all $`0jN`$, $`0<T<t`$, $`T`$ is sufficiently large under the CA (12). Consequently the region $`p_11`$, $`p_21`$ in the DK CA does not show chaos (damage spread).
Proof. Assume that $`x(j,0)=x(j+1,0)=1`$. Then the system (14) implies
$$x(j1,1)=x(j,1)=x(j+1,1)=x(j+2,1)=1.$$
Continue one gets after $`t`$ time steps $`x(k,t)=1`$, where $`jtkj+t+1`$. This proves the first part. Now since the CA (14) will tend to $`x(j,t)=1`$ for all $`0jN`$, $`0<T<t`$, then any change in the initial conditions that preserves the condition $`x(j,0)=x(j+1,0)=1`$ for some $`j`$ will not affect the asymptotic behavior of the CA (14). This completes the proof.
The case $`p_10`$ in the DK model corresponds to the CA
$$x(j,t+1)=x(j1,t)x(j+1,t)\mathrm{mod}\mathrm{\hspace{0.33em}2}.$$
(15)
Following similar steps as those in proposition (3) one can prove the following:
Proposition (4): Any finite initial configuration with two consecutive zeros will tend to the homogeneous configuration $`x(j,t)=0`$ for all $`0jN`$, $`0<T<t`$, $`T`$ is sufficiently large under the CA (15). Consequently the region $`p_11,p_21`$ in the DK CA does not show chaos (damage spread) or periodic configurations.
In the limit $`p_20,p_11`$, DK model corresponds to rule 90
$$x(j,t+1)=x(j1,t)+x(j+1,t)\mathrm{mod}\mathrm{\hspace{0.33em}2},$$
which is known to be chaotic.
All of the above results agree with numerical simulations. Bagnoli et al model \[Bagnoli et al 2002\] is given by
$`\mathrm{If}x(j1,t)+x(j,t)+x(j+1,t)`$ $`=`$ $`0\mathrm{then}x(j,t+1)=0.`$
$`\mathrm{If}x(j1,t)+x(j,t)+x(j+1,t)`$ $`=`$ $`1,`$
$`\mathrm{then}x(j,t+1)`$ $`=`$ $`1\mathrm{with}\mathrm{probability}p_1.`$
$`\mathrm{If}x(j1,t)+x(j,t)+x(j+1,t)`$ $`=`$ $`2,`$ (16)
$`\mathrm{then}x(j,t+1)`$ $`=`$ $`1\mathrm{with}\mathrm{probability}p_2.`$
$`\mathrm{If}x(j1,t)+x(j,t)+x(j+1,t)`$ $`=`$ $`3\mathrm{then}x(j,t+1)=1.`$
where $`x(j,t)`$ are Boolean variables. The limit $`p_11,p_21`$ corresponds to the CA
$`x(j,t+1)`$ $`=`$ $`x(j1,t)+x(j+1,t)+x(j,t)+x(j1,t)x(j+1,t)+`$
$`x(j,t)x(j+1,t)+x(j1,t)x(j,t)+`$
$`x(j1,t)x(j,t)x(j+1,t)\mathrm{mod}\mathrm{\hspace{0.33em}2}.`$
The limit $`p_10,p_20`$ corresponds to the CA
$$x(j,t+1)=x(j1,t)x(j,t)x(j+1,t)\mathrm{mod}\mathrm{\hspace{0.33em}2}.$$
(18)
Proposition (5): a) Any nonzero finite initial configuration will evolve under the CA (17) into the homogeneous configuration $`x(j,t)=1`$ for all $`0jN`$, for $`t`$ is sufficiently large. Hence the limit $`p_11,p_21`$ in Bagnoli et al model does not show chaos or periodic configurations.
b) Any finite initial configuration containing at least one zero site will evolve under the CA (18) into the homogeneous configuration $`x(j,t)=0`$ for all $`0jN`$, for $`t`$ is sufficiently large. Hence the limit $`p_10,p_20`$ in Bagnoli et al model does not show chaos or periodic configurations.
Proof. similar to proposition (3).
The limit $`p_11,p_20`$ corresponds to the CA
$`x(j,t+1)`$ $`=`$ $`x(j1,t)+x(j+1,t)+x(j,t)+`$ (19)
$`2x(j1,t)x(j,t)x(j+1,t)\mathrm{mod}\mathrm{\hspace{0.33em}2}`$
which is similar to rule 150 $`x(j,t+1)=x(j1,t)+x(j+1,t)+x(j,t)`$, hence chaos is expected in Bagnoli et al model in this limit. All of the above results agree with numerical simulations.
It is interesting how CA unite polynomials on finite fields, circulant matrices, graph theory techniques and many other branches of mathematics into one branch which is important both mathematically and from the point of view of applications in complex systems.
## Appendix (2): Overview of networks in CAS
Complex systems are often modeled as graphs where agents are the vertices and the interactions form the edges of the graph. Typically graphs are either regular lattices (e.g. square or cubic), random or scale free where the probability that a vertex has degree $`k`$ is $`p(k)k^\gamma `$. Most of the real networks are of the scale free type. Some proposed mechanisms for this fat tailed distribution \[Dorogovtsev and Mendes 2004\] are self organization (c.f. biological systems) and optimization involving many agents (c.f. economy).
Random graphs were first studied by the mathematicians Erdös and Rényi \[Erdös and Rényi 1960\]. Their model consists of $`N`$ nodes, such that every pair of nodes is connected by a bond with probability $`p`$. The recent increase in computing power and the appearance of interdisciplinary sciences has lead to a better understanding of the properties of complex networks.
Two main properties of complex networks are clustering and small world effect.
Small-world effect means the average shortest node to node (vertex to vertex) distance is very short compared with the whole size of the system (total number of vertices). For social networks, the social psychologist Milgram \[Milgram 1967\] concluded that the average length of the path of acquaintances connecting each pair of people in the United States is six. This concept is known as the six degrees of separation. Such an effect makes it easier for an effect (e.g. an epidemic) to spread throughout the network.
In a regular 1-dimensional lattice of size $`N`$, the average shortest path connecting any two vertices $`l`$ increases linearly with the system size. So regular lattices do not display small-world effect. On the other hand for a random graph, with coordination number $`z`$, one has $`z`$ first (nearest) neighbors, $`z^2`$ second neighbors and so on. This means that the total number of vertices $`N=z^l`$, this gives
$$l=\frac{\mathrm{Ln}(N)}{\mathrm{Ln}(z)}.$$
The logarithmic increase with the size of the lattice allows the distance $`l`$ to be very short even for large $`N`$. Then random graphs display the small-world effect.
Clustering is a common property of complex networks. It means that every vertex has a group of connected nearest neighbours (NN) (collaborators, friends), some of them will often be a connected NN to another vertex. As a measure for the clustering property, a clustering coefficient $`C`$ is defined as the probability that connected pairs of NN of a vertex are also connected to each others. For a random graph, $`C=z/N`$ which goes to zero for large $`N`$. So random graphs do not display clustering property. On the other hand, a fully connected regular lattice itself forms a cluster, then its cluster coefficient is equal to $`1`$.
Complex networks display a small-world effect like random graphs, and they have large clustering coefficient as regular lattices. For a review on many real-world examples, see \[Dorogovtsev and Mendes 2004\].
A small-world network (SWN) proposed initially by Watts and Strogatz \[Watts and Strogatz 1998\] is a superposition of a regular lattice (with high clustering coefficient) and a random graph (with the small world effect). SWN satisfy the main properties of social networks. Also, the structure of SWN combines between both local and nonlocal interactions which is observed in many real systems. For example epidemic spreading show nonlocal interactions e.g SARS.
The concept of SWN has been applied successfully in modelling many CAS, e.g, some games \[Ahmed and Elgazzar 2000 a\], epidemics \[Ahmed et. al. 2002\], economic systems \[Elgazzar 2002\], and opinion dynamics \[Elgazzar 2001\].
An important property related to disease spread in a network is the second moment of the degree distribution i.e. $`k^2`$. If it is divergent then on average a vertex has an infinite number of second nearest neighbors thus if a single vertex is infected the disease will spread in the whole network. This explains the results that disease spread on scale free networks has zero threshold (contrary to the ODE and PDE models). However one should realize that real networks are finite hence a kind of threshold is expected.
Scale-free networks \[Albert and Barabási 2002\] are another class of complex networks. A scale-free network does not have a certain scale. Some nodes have a huge number of connections to other nodes, whereas most nodes have only a few, following a power law distribution.
## Appendix (3): Basics of game theory
Game theory \[Hofbauer and Sigmund 1998\] is the study of the ways in which strategic interactions among rational players produce outcomes (profits) with respect to the preferences of the players. Each player in a game faces a choice among two or more possible strategies. A strategy is a predetermined program of play that tells the player what actions to take in response to every possible strategy other players may use. A basic property of game theory is that one’s payoff depends on the others’ decisions as well as his.
The mathematical framework of the game theory was initiated by von Neumann and Morgenstern in 1944. Also they had suggested the max-min solution for games which is calculated as follows: Consider two players A and B are playing against each other. Two strategies $`S_1`$, $`S_2`$ are allowed for both of them. This game is called two-player, two-strategy game. Assume that the constants $`a,b,c`$ and $`d`$ represent the payoffs (profits) such that, if the two players use the same strategy $`S_1(S_2)`$, their payoff is $`a(d)`$. When a player with strategy $`S_1`$ plays against another one with strategy $`S_2`$, the payoff of the $`S_1`$-player is $`b`$ and the payoff of the $`S_2`$-player is $`c`$ and so on. This is summarized in the payoff matrix as follows:
$$\begin{array}{ccc}& S_1& S_2\\ S_1& a& b\\ S_2& c& d\end{array}.$$
The max-min solution of von Neumann and Morgenstern is for the first player to choose max{min$`(a,b)`$,min$`(c,d)`$}. The second player chooses min{max$`(a,c)`$, max$`(b,d)`$}. If both quantities are equal then the game is stable. Otherwise use mixed strategies.
A weakness of this formalism has been pointed out by Maynard Smith in the hawk-dove (HD) game whose payoff matrix is
$$\mathrm{\Pi }=\begin{array}{ccc}& \mathrm{H}& \mathrm{D}\\ \mathrm{H}& \frac{1}{2}(vc)& v\\ \mathrm{D}& 0& \frac{v}{2}\end{array}.$$
The max-min solution implies (for $`v<c`$) that the solution is D yet as he pointed out this solution is unstable since if one of the players adopts H in a population of D he will have a very large payoff which will make other players switch to H and so on till number of H is large enough that they play each other frequently and get the low payoff $`(vc)/2`$. Thus the stable solution is that the fraction of hawks should be nonzero. To quantify this concept one may use the replicator equation which intuitively means that the rate of change of the fraction of players adopting strategy $`i`$ is proportional to the difference between their payoff and the average payoff of the population i.e.
$$\frac{\mathrm{d}x_i}{\mathrm{d}t}=x_i\left[(\mathrm{\Pi }x)_ix\mathrm{\Pi }x\right],i=1,2,\mathrm{},n,\underset{i=1}{\overset{n}{}}x_i=1,$$
(20)
where $`x_i`$ is the fraction of players adopting strategy $`i`$, and $`\mathrm{\Pi }`$ is the payoff matrix. Applying Eq. (20) to the HD game, one gets that the asymptotically stable equilibrium solution is $`x=v/c`$, where $`x`$ is the fraction of hawks in the population.
For asymmetric game the replicator dynamics equation is
$$\frac{\mathrm{d}x_i}{\mathrm{d}t}=x_i\left[(\mathrm{\Pi }_1y)_ix\mathrm{\Pi }_1y\right],\frac{\mathrm{d}y_i}{\mathrm{d}t}=y_i\left[(\mathrm{\Pi }_2x)_iy\mathrm{\Pi }_2x\right],i=1,2,\mathrm{},n.$$
A basic drawback of normal game theory is the assumption that all players interact globally. It is more realistic to study local games \[Ahmed and Elgazzar 2000 b\] e.g. games on a lattice where players interact only with their nearest neighbors. Also there are several modifications for game formulations.
## Appendix (4): Unpredictability in CAS
There are at least two sources for unpredictability in CAS. The first is that CAS are open systems hence perturbing a CAS may affect another related one e.g. the insect population affected by the perturbation of Lake Victoria. Another reason is the nonlinear interactions \[Strogatz 2000\] between the elements of the CAS. The scientific and mathematical study of Chaos Theory contains many overlaps with the study of Complex Systems, but with differences related to method: Chaos Theory can be used to study Complex Systems, but is not restricted to the study of these systems. Chaos Theory ”deals with deterministic systems whose trajectories diverge exponentially over time” (Bar Yam, NECSI website). It has been used to study Complex Systems, because these systems can be generally defined as a ”deterministic system that is difficult to predict”. On the other hand, complexity deals with systems composed of many interacting agents” The point being that Chaos Theory is one of many tools and methods that can be applied to the study of Complex Systems, but is not specifically devoted to the way these systems are designed, developed, studied, and modeled. That being stated, the famous example of the ”Butterfly Effect” in a chaotic system is an example of an agent (a butterfly) evoking a non-linear response (the storm in New England) within a Complex System (Global Weather System).
A simple example of nonlinear interactions is the logistic difference equation
$$x_{t+1}=rx_t(1x_t),t=0,1,2,\mathrm{},n,r>0.$$
(21)
This equation has two equilibrium solutions $`x=0,x=11/r(r>1)`$ which are asymptotically stable if $`r<1`$ or $`1<r<3`$ respectively. If $`3<r<3.6`$ then cycles appears and if $`r>3.6`$ chaos sets in. Intuitively chaos is sensitive dependence on initial conditions (for more mathematical definition see \[Holmgren 1996\]). Hence in chaotic systems one cannot make long range predictions c.f. weather. A useful measure of chaos are Lyapunov exponents
$$\lambda =\frac{1}{n}\underset{t=0}{\overset{n1}{}}\mathrm{Ln}\left|\stackrel{´}{f}(x_t)\right|.$$
(22)
Since CAS consists of several interacting agents one studies coupled systems e.g. coupled map lattices \[Kaneko 1993\] given by
$$x_i^{t+1}=(1D)f(x_i^t)+\frac{D}{2}\left[f(x_{i1}^t)+f(x_{i+1}^t)\right],i=1,2,\mathrm{},n.$$
(23)
The homogeneous equilibrium is given by $`x=f(x)`$ and it is asymptotically stable if \[Ahmed and Hegazi 2002\]
$$\left|\stackrel{´}{f}(x)\left[(1D)+D\mathrm{cos}(\frac{k\pi }{n})\right]\right|<1,k=0,1,\mathrm{},n1.$$
(24)
The more realistic case is to assume that the map depends on the agents e.g.
$$x_i^{t+1}=(1D)f_i(x_i^t)+\frac{D}{2}\left[f_{i1}(x_{i1}^t)+f_{i+1}(x_{i+1}^t)\right],i=1,2,\mathrm{},n.$$
(25)
But analytic studies for Eq. (25) are more difficult.
These systems shed some light on how to control (synchronize) some CAS \[Ahmed et al 2003\]. One may increase the coupling constant $`D`$. Also if the network of the agents is more connected (e.g. SWN), then the system is easier to synchronize. Finally external control can be applied preferably at highly connected sites.
## Appendix (5): Elements of multi-objective optimization
Almost every real life problem is multi-objective (MOB) \[Collette and Siarry 2003\]. Methods for MOB optimization are mostly intuitive.
Definition (5): A MOB problem is:
$$\mathrm{Minimize}(\mathrm{min})Z_i(\underset{¯}{x}),i=1,2,\mathrm{},k,\mathrm{subject}\mathrm{to}\underset{¯}{g}(\underset{¯}{x})0,\underset{¯}{h}(\underset{¯}{x})0.$$
(26)
Definition (6): A vector $`\underset{¯}{x}^{}`$ dominates $`\underset{¯}{\overset{´}{x}}`$ if $`Z_i(\underset{¯}{x})Z_i(\underset{¯}{\overset{´}{x}})i=1,2,\mathrm{},k`$ with strict inequality for at least one $`i`$, given that all constraints are satisfied for both vectors.
A non-dominated solution $`\underset{¯}{x}^{}`$ is called Pareto optimal and the corresponding vector $`Z_i(\underset{¯}{x}^{}),i=1,2,\mathrm{},k`$ is called efficient. The set of such solutions is called a Pareto set.
Now we discuss some methods for solving MOB problems:
The first method is the lexicographic method. In this method objectives are ordered according to their importance. Then a single objective problem is solved while completing the problem gradually with constraints i.e.
$$\begin{array}{c}\mathrm{min}Z_1\mathrm{subject}\mathrm{to}\hfill \\ \underset{¯}{g}(\underset{¯}{x})0,\underset{¯}{h}(\underset{¯}{x})=0\hfill \end{array},$$
(27)
then if $`\mathrm{ZMIN}(1)`$ is the solution, the second step is $`\mathrm{min}Z_2`$ subject to $`Z_1=\mathrm{ZMIN}(1)`$, and the constraints in Eq. (26), and so on.
A famous application is in university admittance where students with highest grades are allowed in any college they choose. The second best group is allowed only the remaining places and so on. This method is useful but in some cases it is not applicable.
Proposition (6): An optimal solution for the lexicographic problem is Pareto optimal.
Proof. Let $`\underset{¯}{x}^{}`$ be the solution to the Lexicographic problem $`P_l`$. Thus
$$\underset{¯}{x}\underset{¯}{x}^{},\mathrm{then}Z_i(\underset{¯}{x})=Z_i(\underset{¯}{x}^{}),i=1,2,\mathrm{},l1\mathrm{and}Z_l(\underset{¯}{x}^{})<Z_l(\underset{¯}{x}).$$
(28)
Thus $`\underset{¯}{x}^{}`$ is not dominated.
The second method is the method of weights. Assume that it is required to minimize the objectives $`Z(j),j=1,2,\mathrm{},n`$. (The problem of maximization is obtained via replacing $`Z(j)`$ by $`Z(j)`$. Define
$$Z=\underset{i=1}{\overset{k}{}}Z_iw(i),\mathrm{\hspace{0.33em}0}w(i)1,\underset{i=1}{\overset{k}{}}w(i)=1.$$
(29)
Then the problem becomes to minimize $`Z`$ subject to the constraints. This method is easy to implement but it has several weaknesses. The first is that it is not applicable if the feasible set is not convex. The second difficulty of this method is that it is difficult to apply for large number of objectives. However it is quite effective for multiobjective problems with discrete parameters since in this case Pareto optimal set is discrete not a continuous curve.
The third method is the compromise method (sometimes called $`\epsilon `$constr-aint method $`P_\epsilon (k)`$. In this case one minimizes only one objective while setting the other objectives as constraints e.g. minimize $`Z(k)`$ subject to $`Z(j)a(j),j=2,3,\mathrm{},k1,k+1,\mathrm{},n`$, where $`a(j)`$ are parameters to be gradually decreased till no solution is found. The problem with this method is the choice of the thresholds $`a(j)`$. If the solution is unique, then this method is guaranteed to give a Pareto optimal solution.
Proposition (7): If the solution is unique, then the $`\epsilon `$constraint method is guaranteed to give a Pareto optimal solution.
Proof. Let $`\underset{¯}{x}^{}`$ be the optimal solution for the $`\epsilon `$constraint method then
$$\underset{¯}{x}\underset{¯}{x}^{},\mathrm{then}Z_k(\underset{¯}{x}^{})<Z_k(\underset{¯}{x}),$$
hence $`\underset{¯}{x}^{}`$ is Pareto optimal. If $`\underset{¯}{x}^{}`$ is not unique, then it is weakly Pareto i.e. there is no $`\underset{¯}{x}\underset{¯}{x}^{}`$ such that $`Z_i(\underset{¯}{x}^{})<Z_i(\underset{¯}{x})i=1,2,\mathrm{},n`$.
A fourth method using fuzzy logic is to study each objective individually and find its maximum and minimum say $`\mathrm{ZMAX}(j)`$, $`\mathrm{ZMIN}(j)`$, respectively. Then determine a membership $`m(j)=(\mathrm{ZMAX}(j)Z(j))/(\mathrm{ZMAX}(j)\mathrm{ZMIN}(j))`$. Thus $`0m(j)1`$. Then apply $`\mathrm{max}\{\mathrm{min}\{m(j),j=1,2,,n\}\}`$. Again this method is guaranteed to give a Pareto optimal solution provided that the solution is unique otherwise it is weakly Pareto. This method is a bit difficult to apply for large number of objectives. A fifth method is Keeney-Raiffa method which uses the product of objective functions to build an equivalent single objective one.
## Appendix (6): Fractional calculus in CAS
Recently \[Stanislavsky 2000\] it became apparent that fractional equations solve some of the above mentioned problems for the PDE approach. To see this consider the following evolution equation
$$\frac{\mathrm{d}f(t)}{\mathrm{d}t}=\lambda ^2_0^tk(t\stackrel{´}{t})f(\stackrel{´}{t})d\stackrel{´}{t}.$$
(30)
If the system has no memory then $`k(t\stackrel{´}{t})=\delta (t\stackrel{´}{t})`$ and one gets $`f(t)=f_0\mathrm{exp}(\lambda ^2t)`$. If the system has an ideal memory, then
$$k(t\stackrel{´}{t})=\{\begin{array}{c}1,t\stackrel{´}{t}\hfill \\ 0,t<\stackrel{´}{t}\hfill \end{array},$$
hence $`ff_0\mathrm{cos}(\lambda t)`$. Using Laplace transform
$$L[f]=_0^{\mathrm{}}f(t)\mathrm{exp}(st)dt,$$
one gets $`L[f]=1`$ if there is no memory and $`L[f]=1/s`$ if there is ideal memory hence the case of non-ideal memory is expected to be given by $`L[f]=1/s^\alpha ,\mathrm{\hspace{0.33em}0}<\alpha <1`$. In this case Eq. (28) becomes
$$\frac{\mathrm{d}f(t)}{\mathrm{d}t}=_0^t\frac{(t\stackrel{´}{t})^{\alpha 1}f(\stackrel{´}{t})\mathrm{d}\stackrel{´}{t}}{\mathrm{\Gamma }(\alpha )},$$
(31)
where $`\mathrm{\Gamma }(\alpha )`$ is the Gamma function. This system has the following solution
$$f(t)=f_0E_{\alpha +1}(\lambda ^2t^{\alpha +1}),$$
where $`E_\alpha (z)`$ is the Mittag Leffler function given by
$$E_\alpha (z)=\underset{k=0}{\overset{\mathrm{}}{}}\frac{z^k}{\mathrm{\Gamma }(\alpha k+1)}.$$
It is direct to see that $`E_1(z)=\mathrm{exp}(z),E_2(z)=\mathrm{cos}(z)`$
Following a similar procedure to study a random process with memory, one obtains the following fractional evolution equation
$$\frac{^{\alpha +1}P(x,t)}{t^{\alpha +1}}=\underset{n}{}\frac{(1)^n}{n!}\frac{^n[K_n(x)P(x,t)]}{x^n},\mathrm{\hspace{0.33em}0}<\alpha <1,$$
(32)
where $`P(x,t)`$ is a measure of the probability to find a particle at time $`t`$ at position $`x`$.
We expect that Eq. (30) will be relevant to many complex adaptive systems and to systems where fractal structures are relevant since it is argued that there is a relevance between fractals and fractional differentiation \[Rocco and West 1999\].
For the case of fractional diffusion equation the results are
$`{\displaystyle \frac{^{\alpha +1}P(x,t)}{t^{\alpha +1}}}`$ $`=`$ $`D{\displaystyle \frac{^2P(x,t)}{x^2}},P(x,0)=\delta (x),{\displaystyle \frac{P(x,0)}{t}}=0`$
$`P`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{D}t^\beta }}M({\displaystyle \frac{\left|x\right|}{\sqrt{D}t^\beta }};\beta ),\beta ={\displaystyle \frac{\alpha +1}{2}},`$ (33)
$`M(z;\beta )`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^nz^n}{n!\mathrm{\Gamma }(\beta n+1\beta )}}.`$
For the case of no memory $`\alpha =0M(z,1/2)=\mathrm{exp}(z^2/4)`$.
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# A Dark Matter Candidate With New Strong Interactions
## 1 Introduction
The most common particle physics models for dark matter involve weakly interacting particles. They can be broadly classified as WIMPS or axions, with the theoretician’s favorite WIMP being a neutralino of the Supersymmetric Standard Model (SSM). Within string theory, the physics of both of these candidates is closely connected to SUSY breaking, because string theory axions generally arise from moduli fields, whose mass is related to a superpotential on moduli space.
One of the authors has recently introduced a new model for SUSY breaking, which has no candidate for either WIMP or axion dark matter. The model is based on the principle of Cosmological SUSY Breaking (CSB):
* The (positive) cosmological constant (c.c.) is a discrete tunable parameter, governing the number of states in the Hilbert space of quantum gravity in de Sitter (dS) space.
* As the c.c. vanishes, SUSY is restored, with the relation $`m_{3/2}\mathrm{\Lambda }^{1/4}`$ between the gravitino mass and the c.c. A discrete $`Z_n`$ $`R`$ symmetry is restored in the same limit, explaining, in low energy terms, the vanishing of the c.c. in the SUSic limiting theory. The limiting theory must have a compact moduli space, in order to guarantee that the dS state of the low energy effective field theory is stable.
* SUSY breaking is spontaneous in the low energy effective theory, but is induced by $`R`$ breaking terms in the Lagrangian which have no low energy explanation. The coefficients in these terms are tuned to guarantee the CSB scaling relation between $`m_{3/2}`$ and $`\mathrm{\Lambda }`$.
As a consequence of the first requirement, the low energy effective field theory of CSB must contain a goldstino field: a linear supermultiplet which is massless in the SUSic, R symmetric limit. In this was taken to be a chiral superfield $`G`$, with $`R`$ charge $`0`$. If there are no fields of $`R`$ charge $`2modn`$ in the low energy theory, then $`G`$ is naturally massless. $`R`$ charges were assigned to standard model fields in a way that insured the absence of all baryon and lepton number violating dimension $`4`$ and $`5`$ operators, apart from the term $`n^{ij}H_u^2L_iL_j`$ (which gives rise to neutrino masses). The generation of this term, and of the texture of Yukawa couplings is imagined to have to do with physics at the unification scale. There is also an ordinary discrete symmetry $``$, under which $`G`$ transforms. $``$ allows the coupling $`g_\mu GH_uH_d`$ but forbids the conventional $`\mu `$ term. $`G^a`$ is the lowest order $``$ invariant monomial in $`G`$.
High energy physics supplies us with a term $`M_P^2\mathrm{\Lambda }^{1/4}f(G/M_P)`$ which violates $`R`$ and implements CSB. The dimensionless coefficients in the function $`f`$ are tuned to guarantee that the c.c. is indeed $`\mathrm{\Lambda }`$. For phenomenological reasons, one must also add terms
$$d^4\theta M_1^2K(g,h_u,h_d,q,\overline{u},\overline{d},l,\overline{e}),$$
and,
$$d^2\theta Z_A(g^a)W_A^2+h.c..$$
We have used an unconventional notation where a lower case label $`s`$ for a chiral superfield $`S`$ stands for $`S/M_1`$. The Kahler potential depends, of course, both on chiral fields and their conjugates. The functions $`K`$ and $`Z_A`$ are imagined to emerge from integrating out degrees of freedom at a scale $`M_1M_UM_P`$, whose value is determined by RG flow in the limiting $`\mathrm{\Lambda }=0`$, theory. They can be chosen to satisfy all phenomenological requirements if $`M_11`$ TeV. It is easy to invent strongly coupled theories $`𝒢`$ which could give rise to all the required properties save one. There is no known example of a theory which preserves the $`R`$ symmetry, and leaves exactly one effective chiral superfield which could play the role of $`G`$. We will leave this problem to future work and concentrate on the problem of dark matter.
If the coupling functions $`Z_A`$ were forced to be logarithms by an accidental $`U(1)`$ with standard model anomalies, then the real part of $`G`$ could be a QCD axion. However, it would have a range of axion couplings ruled out by beam dump experiments. Consequently the model has no axion candidates. The basic setup of CSB contradicts the idea of SUSY neutralino dark matter. The gravitino is the LSP in the CSB scenario, and its longitudinal components are relatively strongly coupled, so the NLSP is not cosmologically stable.
The only plausible dark matter candidate in this scenario is what we will call a $`𝒢`$ baryon. That is, we assume the strongly interacting $`𝒢`$ sector has an accidental symmetry, which renders the lightest particle carrying some accidental $`U(1)`$ quantum number, cosmologically stable. In this paper, we will explore the idea that the dark matter is in fact a baryon of a strongly interacting sector with an RG scale of order $`M_1`$. We will see that under a variety of assumptions about the production of this particle, this hypothesis is consistent with conventional cosmology. It has the added virtue of correlating the coincidence between the dynamical scale $`M_1`$ and the CSB scale $`\sqrt{(\mathrm{\Lambda }^{1/4}M_P)}`$ to the existence of galaxies. That is to say, we imagine that the limiting model calculates the value of the scale $`M_1`$ and the other parameters of $`e.g.`$ the inflaton field, in such a way that the density of $`𝒢`$ baryons coincides with what we know about dark matter density from observations. Now consider the model of CSB, with various values of $`\mathrm{\Lambda }`$. The only values which will produce a model with galaxies will be those which satisfy Weinberg’s bound. At least within a few orders of magnitude, this matches the scale of CSB to $`M_1`$ and the dark energy density to the dark matter density (cosmic coincidence).
We will also see that there is a variety of thermal histories for the universe in which $`𝒢`$ baryons can be dark matter only if there is a CP violating $`𝒢`$ baryon number asymmetry. We might imagine a model in which $`𝒢`$ and ordinary baryon asymmetries were produced by the same mechanism, perhaps explaining the dark/baryonic matter ratio of the universe.
The Hess telescopes have seen a photon signal from the center of the galaxy, which might be consistent with a dark matter candidate of mass $`1518`$ TeV, if dark matter in the galaxy follows the profile predicted by . It is very hard to find a neutralino model which can produce such a large mass, basically because weak annihilation cross sections decrease with mass. On the other hand, strongly interacting particles have mass independent annihilation cross sections and can easily fit this data.
In the next section, we estimate various cross sections for baryon like objects, using large $`N`$ QCD as a paradigm. The $`𝒢`$ theory must differ from QCD since it preserves chiral symmetry and is supersymmetric. Nonetheless, we hope that these estimates give us a rough guide to the scales involved. We then go on to estimate the mass, cross section and primordial asymmetry for which a $`𝒢`$ baryon could be dark matter. We consider two scenarios: a standard thermal relic abundance calculation, and a particular non-thermal production scheme. We find that for reasonable values of parameters, the model can fit the data, and perhaps reproduce the Hess signal. To answer the latter question in more detail, one must perform a detailed estimate of the photon spectrum one would get from annihilation processes involving a strongly interacting dark matter candidate. We are not sure that the model used by the Hess collaboration in order to extract the parameters of a hypothetical dark matter particle from their signal, takes into account the physics of a strongly interacting particle.
We should emphasize that despite our original motivation, our calculations would be applicable to any dark matter candidate with new strong interactions of the right scale. In particular, we note that our model for dark matter is similar to the hypothesis that dark matter is a techni-baryon.
## 2 Annihilation Cross Sections for Dark Matter With New Strong Interactions
The nucleon anti-nucleon annihilation cross section is usually written in units of the pion Compton wavelength, because this is the range of nuclear forces. In fact, this parametrization is singular in the chiral limit, when the pion becomes a Goldstone boson. It is not correct that the cross section blows up in this limit.
A better estimate is obtained by thinking about chiral soliton models of the nucleon. In such models the nucleon is realized as a classical solution of a large $`N`$ effective action. The effective Planck constant of this action is of order $`N`$, and the scale over which solutions vary is the QCD scale. Although these models use the spontaneously broken chiral symmetry of QCD in an essential way, they give the same order of magnitude results one would expect from general large $`N`$ considerations. We expect the size of a general large $`N`$ soliton to be given by such an $`N`$ independent scale, and large $`N`$ soliton masses will be of order $`N`$.
The soliton-anti-soliton annihilation cross section will be given by its classical size $`\sigma \mathrm{\Lambda }_𝒢^2`$ and will be more or less energy independent in the regime of interest, because the cosmological velocities of these heavy particles will be low. Note that this is not s-wave annihilation. The typical orbital angular momentum involved in these collisions is of order $`\frac{\sqrt{m_𝒢T}}{\mathrm{\Lambda }_𝒢}`$, where $`T`$ is the temperature at which the annihilation takes place. Note also that the thermally averaged cross section $`<\sigma v>`$, which appears in cosmological Boltzmann equations, will be $`O(T/m_G)^{1/2}`$. We believe that this is the correct scaling even for ordinary baryons, and that conventional calculations of the relic baryon density in a baryon symmetric universe are not quite correct. However, this does not change the qualitative conclusion of those calculations, namely that we need a baryon asymmetry to account for the observed baryon number density of the universe.
We note that the reason that we are interested in large $`N`$ counting is the combination of the Hess data, and the constraints on $`\mathrm{\Lambda }_𝒢`$ from supersymmetric phenomenology. The latter prefers a scale $`\mathrm{\Lambda }_𝒢1`$ TeV, in order to accommodate the bounds on charged superpartner masses, while the former indicates a mass around $`1518`$ TeV for the dark matter particle. In a large $`N`$ model, the baryon mass would be $`N\alpha \mathrm{\Lambda }_𝒢`$ with $`\alpha `$ a number of order $`1`$ ($`\alpha 2`$ in QCD). Thus, we would want $`N\alpha 1518`$. These are not unreasonable values. For example, the best of the inadequate models for the $`𝒢`$ theory, studied in was an $`SU(4)`$ SUSY gauge theory. For $`N=4`$, we require $`\alpha 4`$, about twice the value in QCD.
We emphasize however that we do not know the details of the model which the Hess collaboration used in quoting $`1518`$ TeV for their best fit to the dark matter candidate. In particular, for weakly coupled neutral dark matter, the direct photon annihilation signal is suppressed by a power of $`(\alpha /\pi )^2`$ relative to photons produced from decays of particles with direct coupling to the dark matter. There is no such suppression for strongly interacting neutral composites of charged particles. For example the large $`N`$ nucleon magnetic moment is order $`e(=\sqrt{4\pi \alpha _{em}})N`$ in $`\mathrm{\Lambda }_{QCD}`$ units. Hess has not yet seen the characteristic turnover in their photon signal, which would be expected from dark matter annihilation, and the question of astrophysical explanations for the signal from the galactic center is still controversial. It is perhaps premature to try to fit their spectrum.
However, it is clear that in order to really confront an eventual dark matter signal from Hess data, we need a much better estimate of the photon spectrum produced by a $`𝒢`$ baryon. In addition, since we find that for most values of the reheat temperature of the universe, we must invoke a $`𝒢`$ baryon asymmetry to account for the observed dark matter density, the annihilation signal will be proportional to the small density of anti-$`𝒢`$ baryons. We have not yet done the calculations to determine the range of parameters for which we would expect a significant annihilation signal from the center of the galaxy. In the rest of this paper, we will choose an annihilation cross section of order $`\mathrm{\Lambda }_𝒢^2`$ and parametrize our results in terms of the $`𝒢`$ baryon mass $`m_𝒢>\mathrm{\Lambda }_𝒢`$, $`\mathrm{\Lambda }_𝒢`$, and an asymmetry.
Our description of the $`𝒢`$ baryon will utilize the following characteristics of a soliton model: energy independent annihilation cross section much larger than the scale of its Compton wavelength, and thermal production at energies well below its mass. The latter is a well known characteristic of solitons in weakly coupled field theory. Finally, we will parametrize the $`𝒢`$ baryon mass as $`N\alpha \mathrm{\Lambda }_𝒢`$, with $`\mathrm{\Lambda }_𝒢1`$ TeV, in order to suggest the large $`N`$ scaling of soliton masses in strongly coupled gauge theories with large gauge groups.
## 3 The Relic Abundance of $`𝒢`$ baryons
We will denote by $`\mathrm{\Omega }_G`$ the fraction of the observed density of the universe in $`𝒢`$ baryons plus anti-baryons. To match the observed dark matter abundance, we require $`\mathrm{\Omega }_G\frac{\rho _G}{\rho _{cr}}=.24`$, using the data from WMAP which specifies $`\mathrm{\Omega }_m=.29\pm .07`$ and $`\mathrm{\Omega }_b=.047\pm .006`$. If $`n_G`$ is the number of $`𝒢`$ baryons per comoving volume, then this can be written $`\mathrm{\Omega }=\frac{n_Gm_𝒢}{3H_0^2/8\pi G}=\frac{n_Gm_𝒢}{1.054h^210^4\frac{eV}{cm^3}}`$.
Writing today’s value of the $`𝒢`$ baryon abundance (the ratio of the number of $`𝒢`$ baryons per comoving volume and the entropy) $`Y_0\frac{n_G}{s_0}`$, this condition becomes
$$\mathrm{\Omega }=.24=\frac{s_0Y_0m_𝒢}{1.054h^210^4\frac{eV}{cm^3}}$$
Thus we require $`Y_0=\frac{.44eV}{m_𝒢}`$. We will write $`m_𝒢=N\alpha `$ TeV, treating $`1`$ TeV as the analog of the QCD scale for the $`𝒢`$ gauge theory, and applying a large $`N`$ scaling rule for baryon masses. In QCD $`N=3`$ and $`\alpha 2`$. Our point is that the analog of a baryon mass could be quite a bit higher than $`1`$ TeV. For example $`N=6`$ and $`\alpha 3`$ would give us an $`18`$ TeV dark matter candidate, as would be required by the interpretation of Hess data in terms of dark matter annihilation. With this parametrization, the required value of the abundance is $`Y_0\frac{410^{13}}{N\alpha }`$.
The relic abundance of $`𝒢`$ baryons depends on some assumptions about the evolution of the universe at the TeV scale and above. We assume that there was a reheating process which gives rise to a radiation dominated universe at some temperature $`T_{RH}`$. This might be due to primordial inflaton decay, or the later decay of some other massive particle which dominates the energy density before it decays. We call the width of the particle $`\mathrm{\Gamma }_X`$ . If $`T_{RH}>1`$ TeV, the $`𝒢`$ gauge theory is thermalized by X-decay and the post-decay distribution of $`𝒢`$ baryons is given by the thermal ensemble. Note that this is true even when $`m_𝒢T_{RH}`$. In this regime of parameters, the $`𝒢`$ baryon is a thermal relic, and we find that, in the absence of an asymmetry, the relic abundance is too small to explain the observed dark matter density.
For $`T_{RH}<1`$ TeV, $`𝒢`$ baryons are produced non-thermally and we must be a bit more specific about the dynamics. For a weakly coupled $`X`$ particle, $`m_XT_{RH}`$ and we can still have $`X`$ decays into $`𝒢`$ baryons. Suppose first that $`m_Xm_𝒢`$ so that we can treat the $`𝒢`$ baryons as just another massless species. If we assume the couplings to $`𝒢`$ baryons are not suppressed relative to standard model particles we get a branching ratio of order $`10^2`$ into $`𝒢`$ baryons. The decay will be reasonably rapid, so we neglect annihilation processes during the decay period and obtain an initial abundance of
$$Y_010^2\frac{T_{RH}}{m_𝒢}.$$
If the $`𝒢`$ baryon were massless, this ratio would just be the branching ratio $`10^2`$. The additional suppression is our estimate of the number of $`𝒢`$ baryons per photon that result from the thermalization process.
Throughout the interesting range of parameters, the X particle life-time is short enough to neglect annihilation in the calculation above. Now we can evolve the resulting $`𝒢`$ baryon densities according a Boltzmann equation driven only by the annihilation of $`g`$ and $`\overline{g}`$. As in the discussion of solitonic dark matter abundance in Griest and Kamionkowski , the thermally averaged annihilation cross sections have temperature dependence given by : $`<\sigma |v|>=\sigma _0(\frac{T_{RH}}{m_𝒢})^{1/2}`$. Also,we will assume that $`T_{RH}<m_𝒢`$. In this case there can be no process in the Boltzmann equation that creates $`𝒢`$ baryons because it is not energetically favorable. The Boltzmann equation for the evolution of $`𝒢`$ baryons is:
$$\dot{n}_g+3Hn_g=<\sigma |v|>n_g^2$$
Letting $`Yn_g/s`$ and $`xm_𝒢/T`$ we get
$$\frac{dY}{dx}=\frac{x^{1/2}\sigma _0sm_{pl}}{1.67g_{}^{1/2}m_𝒢^2}Y^2$$
Since $`s=\frac{2\pi ^2}{45}g_sT^3`$, we can then write:
$$\frac{dY}{dx}=\frac{kY^2}{x^{5/2}}$$
where $`k\frac{m_𝒢2\pi ^2\sigma _0g_sm_{pl}}{1.67g_{}^{1/2}45}`$.
Here we will assume an average $`g_{}g_s50`$.
Defining an order one parameter $`\beta `$ such that $`\sigma _0=\frac{1}{(\beta TeV)^2}`$, $`k4.510^{15}N\alpha /\beta ^2`$.
The solution to this equation is:
$$Y_{final}=\frac{1}{\frac{1}{Y_i}+\frac{2k}{3}(\frac{1}{x_i^{3/2}}\frac{1}{x_f^{3/2}})}$$
Notice a few properties of this solution. The present day temperature is so low that $`\frac{1}{x_f^{3/2}}0`$. Hence either the $`\frac{1}{Y_i}`$ term or the $`\frac{2k}{3}\frac{1}{x_i^{3/2}}`$ term dominates, depending on $`T_{RH}`$. The $`\frac{1}{Y_i}`$ term dominates for $`T_{RH}<.3N\alpha /\beta ^2`$ MeV. A reheat temperature in this range would be inconsistent with nucleosynthesis, so we can ignore this term. Thus, $`Y_f`$ is determined by:
$$Y_f=\frac{3x_i^{3/2}}{2k}$$
where $`x_i=(m_𝒢/T_{RH})`$.
In a general model where we do not fix the mass of the $`𝒢`$ baryon or the exact cross section, we can get an upper bound on $`T_{RH}`$ from our requirement that $`Y_0=\frac{4.410^{13}}{N\alpha }`$:
$$T_{RH}>.008\beta ^{4/3}N\alpha \text{TeV}$$
For reheat temperatures below this value, the $`𝒢`$ baryons will dominate the universe. Thus we find a small window $`1>\frac{T_{RH}}{TeV}>.008\beta ^{4/3}N\alpha `$, where non-thermal, symmetric $`𝒢`$ baryon production could account for the observed properties of dark matter. In particular, for typical values $`N\alpha 10`$ and $`\beta 1`$, we find that this window has a width of about an order of magnitude. However, this range for the reheat temperature does not conform to our prejudice that $`m_𝒢`$ is substantially larger than $`\mathrm{\Lambda }_𝒢`$. We also note that there was no loss of generality in our assumption that $`m_X>>m_𝒢`$. If this assumption is not valid, then $`T_{RH}`$ is quite low, and $`𝒢`$ baryons would be overproduced as long as $`m_X>m_𝒢`$.
For $`T_{RH}>1`$ TeV, the thermal relic abundance is too small to account for the observed dark matter, but we can remedy this by postulating an asymmetry. The simplest possibility is that the asymmetry is generated directly in the decay of the X particle, in which case we have the standard result that
$$Y_0=ϵ_G\frac{T_{RH}}{m_X},$$
where
$$ϵ_G\underset{\text{f}}{}B_\text{f}\frac{\mathrm{\Gamma }_X(X\text{f})\mathrm{\Gamma }_X(X\text{})}{\mathrm{\Gamma }_X}$$
$`\mathrm{\Gamma }`$ is a decay rate, f and are all possible final states, and $`B_\text{f}`$ is the total $`𝒢`$ baryon number of the final state f. $`ϵ_B`$ is the corresponding asymmetry in ordinary baryon number. In order to match the observed dark matter density and the observed baryon density, we need
$$\frac{ϵ_G}{ϵ_B}\frac{1}{2N\alpha }\times 10^2,$$
and
$$ϵ_B\frac{T_{RH}}{m_X}8.6\times 10^{11}.$$
In our model $`(\frac{T_{RH}}{m_X}\frac{\sqrt{m_Xm_P}}{M})`$ <sup>2</sup><sup>2</sup>2$`M`$ is the scale of irrelevant couplings of the $`X`$ particle to the standard model, which are responsible for the decay. is bounded from below by the requirements that the X couplings to ordinary matter are at most Planck suppressed, and that the $`X`$ is massive enough to produce the $`𝒢`$ baryon in its decays. Thus $`\frac{T_{RH}}{m_X}>(\frac{N\alpha }{2})^{1/2}\times 10^3`$.
We see that the $`ϵ`$ parameters must be very small in order to account for the observed asymmetries. In fact, small baryon number violating branching ratios arise naturally if we assume that X is a slow roll inflaton, $`I`$, with a “natural” potential of the form $`\mu ^4f(I/m_P)`$. A theorem of Nanopoulos and Weinberg tells us that asymmetries can arise only at second order in baryon violating couplings. Let us assume the decay of the inflaton is mediated primarily via dimension $`5`$ operators. Then even if the dimension five couplings involve CP violation and baryon number violation, we will find that $`ϵ_B(\frac{m_I}{m_P})^2`$. We also have the order of magnitude estimate $`\frac{T_{RH}}{m_I}(\frac{m_I}{m_P})^{1/2}`$, so that
$$Y_B(\frac{m_I}{m_P})^{5/2}.$$
This will fit the observed baryon asymmetry if
$$m_I10^4m_P$$
Note that this gives an inflation scale $`\mu `$ close to the unification scale.
In this context we might attempt to explain the further suppression $`\frac{ϵ_G}{ϵ_B}10^3`$ by postulating that (perhaps as a consequence of the R symmetry introduced in ) the leading contribution to the $`𝒢`$-baryon asymmetry comes from the interference of a dimension $`5`$ and dimension $`6`$ coupling of the inflaton, and is suppressed by a further power of $`\frac{m_I}{m_P}`$. This is off by a factor of $`10`$ but our estimates are so crude that we can consider this a success.
Indeed, we proposed this simple model not because we think it has to be right, but to show that reasonable calculations of both the dark matter and baryon abundances can be obtained for our new form of dark matter.
To summarize, we probably need asymmetric production of $`𝒢`$ baryons to make them an acceptable dark matter candidate. We outlined a plausible model of asymmetric production in inflaton decay, which could naturally explain both the baryon asymmetry of the universe and the dark matter density.
## 4 Conclusions
We have shown that a baryon-like state of the new, strongly interacting, $`𝒢`$ theory, which was introduced in to implement Cosmological SUSY breaking, is a promising dark matter candidate. The new strong interaction scale is around $`1`$ TeV and the $`𝒢`$ baryon mass is somewhat higher, perhaps as high as the $`1518`$ TeV needed to fit the Hess data on photons from the center of the galaxy, if the explanation for that data turns out to be dark matter annihilation. We saw that this sort of baryon to interaction scale ratio was natural in the context of large $`N`$ scaling with $`N57`$. While the $`𝒢`$ theory is probably not as simple as $`SU(N)`$ QCD, there is reason to believe that reasonably large baryon masses are a more general phenomenon.
This sort of dark matter candidate allows one to contemplate a simple explanation of the dark matter to baryon ratio, since the asymmetries in baryon and $`𝒢`$ baryon number might have the same physical origin. We need to explain a factor of order $`10^3`$ in these asymmetries, in order to fit the data. We constructed a plausible model in which both asymmetries are generated in inflaton decay. To explain the size of the asymmetries we invoked the R symmetry of and dimensional analysis. There are more scenarios for baryogenesis in the literature than there are authors on this paper, and it is entirely plausible to us that a more elegant mechanism could be found. However, our simple model might work, and it might be the right answer.
Much more work needs to be done to sort out signatures of such hyper-strongly interacting dark matter, as well as to explore a variety of models for the production of baryon and $`𝒢`$ baryon asymmetries. In addition, it will be necessary to find out more about the dynamics of the as yet mysterious $`𝒢`$ theory, which gives rise to these new particles.
## 5 Acknowledgments
We have benefitted from conversations with J. Primack, P.J. Fox, M. Dine, and S. Dimopoulos.
The research of the authors was supported in part by DOE grant number DE-FG03-92ER40689.
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# Can one see the fundamental frequency of a drum?
## 1 Main result
Let us consider an open set $`\mathrm{\Omega }\text{}^n`$ and denote the bottom of the spectrum of its minus Dirichlet Laplacian $`(\mathrm{\Delta })_{\mathrm{Dir}}`$ by $`\lambda (\mathrm{\Omega })`$. (We understand the minus Dirichlet Laplacian as the self-adjoint operator which is the Friedrichs extension of the operator $`\mathrm{\Delta }`$ defined on $`C_0^{\mathrm{}}(\mathrm{\Omega })`$. ) In case when $`\mathrm{\Omega }`$ is a bounded domain with a sufficiently regular boundary, $`\lambda (\mathrm{\Omega })`$ is the lowest eigenvalue of $`\mathrm{\Delta }`$ with the Dirichlet boundary condition on $`\mathrm{\Omega }`$. In the general case we can write
(1.1)
$$\lambda (\mathrm{\Omega })=\underset{uC_0^{\mathrm{}}(\mathrm{\Omega })}{inf}\frac{_\mathrm{\Omega }|u|^2𝑑x}{_\mathrm{\Omega }|u|^2𝑑x}.$$
It follows that $`\mathrm{\Omega }^{}\mathrm{\Omega }`$ implies $`\lambda (\mathrm{\Omega })\lambda (\mathrm{\Omega }^{}))`$. In particular, if $`B_r`$ is an open ball of radius $`r`$, such that $`B_r\mathrm{\Omega }`$, then $`\lambda (\mathrm{\Omega })\lambda (B_r)=C_nr^2`$ where $`C_n=\lambda (B_1)`$. It follows that for the interior radius of $`\mathrm{\Omega }`$, which is defined as
$$r_\mathrm{\Omega }=sup\{r|B_r\mathrm{\Omega }\},$$
we have
$$\lambda (\mathrm{\Omega })C_nr_\mathrm{\Omega }^2.$$
But this estimate is not good for unbounded domains or domains with complicated boundaries. In particular, a similar estimate from below does not hold in general.
The way to improve this estimate is to relax the requirement for $`B_r`$ to be completely inside $`\mathrm{\Omega }`$ by allowing some part of $`B_r`$, which has a “small” Wiener capacity, to stick out of $`\mathrm{\Omega }`$. Namely, let us take an arbitrary $`\gamma (0,1)`$ and call a compact set $`F\overline{B}_r`$ negligible (or, more precisely, $`\gamma `$-negligible) if
(1.2)
$$\mathrm{cap}(F)\gamma \mathrm{cap}(\overline{B}_r).$$
(Here $`\mathrm{cap}(F)`$ denotes the Wiener capacity of $`F`$, $`\overline{B}_r`$ is the closure of $`B_r`$.)
Now denote
$$r_{\mathrm{\Omega },\gamma }=sup\{r|B_r,\overline{B}_r\mathrm{\Omega }\text{ is }\gamma \text{-negligible}\}.$$
This is the interior capacitary radius.
###### Theorem 1.1
Let us fix $`\gamma (0,1)`$. Then there exist $`c=c(\gamma ,n)>0`$ and $`C=C(\gamma ,n)>0`$, such that for every open set $`\mathrm{\Omega }\text{}^n`$
(1.3)
$$cr_{\mathrm{\Omega },\gamma }^2\lambda (\mathrm{\Omega })Cr_{\mathrm{\Omega },\gamma }^2.$$
Explicit values of constants $`c=c(\gamma ,n)`$ and $`C=C(\gamma ,n)`$ are provided below in (3.19) and (4.16) respectively.
Let us formulate some corollaries of this theorem.
###### Corollary 1.2
$`\lambda (\mathrm{\Omega })>0\text{ if and only if }r_{\mathrm{\Omega },\gamma }<\mathrm{}`$.
This corollary gives a necessary and sufficient condition of strict positivity of the operator $`(\mathrm{\Delta })_{\mathrm{Dir}}`$ in $`\mathrm{\Omega }`$.
Since the condition $`\lambda (\mathrm{\Omega })>0`$ does not contain $`\gamma `$, we immediately obtain
###### Corollary 1.3
Conditions $`r_{\mathrm{\Omega },\gamma }<\mathrm{}`$, taken for different $`\gamma `$’s, are equivalent.
Denoting $`F=\text{}^n\mathrm{\Omega }`$ (which can be an arbitrary closed subset in $`\text{}^n`$), we obtain from the previous corollary (comparing $`\gamma =0.01`$ and $`\gamma =0.99`$):
###### Corollary 1.4
Let $`F`$ be a closed subset in $`\text{}^n`$, which has the following property: there exists $`r>0`$ such that
$$\mathrm{cap}(F\overline{B}_r)0.01\mathrm{cap}(\overline{B}_r)$$
for all $`B_r`$. Then there exists $`r_1>0`$ such that
$$\mathrm{cap}(F\overline{B}_{r_1})0.99\mathrm{cap}(\overline{B}_{r_1})$$
for all $`B_{r_1}`$.
This is a new property of capacity which is proved by spectral theory arguments.
Once upon a time Marc Kac formulated a fascinating question: “Can one hear the shape of a drum?” The precise meaning of this question is as follows: is it possible to reconstruct the drum (a bounded domain in $`\text{}^2`$) up to an isometry by the spectrum of its Dirichlet Laplacian?
Theorem 1.1 suggests formulation of a question, which is roughly inverse to the question above: “Can one see the fundamental frequency of a drum?” More precisely, can one find a simple visual image related to a domain in $`\text{}^2`$ (or $`\text{}^n`$), such that it allows to recover the lowest eigenvalue of the Dirichlet Laplacian in this domain, or at least give reasonably good estimates of this eigenvalue? Assuming that our eye (possibly armed by a visual aid device) can filter out the sets of “small” capacity, a partial answer to this question is given by Theorem 1.1.
The inequalities (1.3) for sufficiently small $`\gamma >0`$ (comparable with $`(4n)^{4n}`$) were established in (see also Chapters 10 and 11 in ). Theorem 1.1 provides a substantial improvement, in particular allowing corollaries 1.3 and 1.4 and providing explicit values of the constants.
The proof of Theorem 1.1 is based on the ideas of our paper . Necessary definitions and results about the Wiener capacity can be found e.g. in .
Acknowledgments. We are grateful to Egon Schulte who communicated to us the coverings multiplicity estimate (3.18), and also to Dan Grieser and Elliott Lieb for useful references.
## 2 Preliminaries on capacity
In this section we will recall some definitions and introduce necessary notations. For simplicity we will always assume that $`n3`$.
For every subset $`𝒟\text{}^n`$ denote by $`\mathrm{Lip}(𝒟)`$ the space of (real-valued) functions satisfying the uniform Lipschitz condition in $`𝒟`$, and by $`\mathrm{Lip}_c(𝒟)`$ the subspace in $`\mathrm{Lip}(𝒟)`$ of all functions with compact support in $`𝒟`$ (this will be only used when $`𝒟`$ is open). By $`\mathrm{Lip}_{loc}(𝒟)`$ we will denote the set of functions on (an open set) $`𝒟`$ which are Lipschitz on any compact subset $`K𝒟`$. Note that $`\mathrm{Lip}(𝒟)=\mathrm{Lip}(\overline{𝒟})`$ for any bounded $`𝒟`$.
If $`F`$ is a compact subset in $`\text{}^n`$, then the Wiener capacity of $`F`$ is defined as
(2.1)
$$\mathrm{cap}(F)=inf\left\{_\text{}^n|u(x)|^2𝑑x|u\mathrm{Lip}_c(\text{}^n),u|_F=1\right\}.$$
Note that the infimum does not change if we restrict ourselves to the functions $`u\mathrm{Lip}_c(\text{}^n)`$ such that $`0u1`$ everywhere (see e.g. , Sect. 2.2.1).
We will also need another (equivalent) definition of the Wiener capacity $`\mathrm{cap}(F)`$ for a compact set $`F\overline{B}_r`$. For $`n3`$ it is as follows:
(2.2)
$$\mathrm{cap}(F)=sup\{\mu (F)|_F(xy)𝑑\mu (y)1\text{on}\text{}^nF\},$$
where the supremum is taken over all positive finite Radon measures $`\mu `$ on $`F`$ and $`=_n`$ is the standard fundamental solution of $`\mathrm{\Delta }`$ in $`\text{}^n`$ i.e.
$$(x)=\frac{1}{(n2)\omega _n}|x|^{2n},$$
with $`\omega _n`$ being the area of the unit sphere $`S^{n1}\text{}^n`$. The maximizing measure in (2.2) exists and is unique. We will denote it $`\mu _F`$ and call it the equilibrium measure. Note that
(2.3)
$$\mathrm{cap}(F)=\mu _F(F)=\mu _F(\text{}^n)=\mu _F,1=_F𝑑\mu _F.$$
The corresponding potential will be denoted $`P_F`$, so
$$P_F(x)=_F(xy)𝑑\mu _F(y),x\text{}^nF.$$
We will call $`P_F`$ the equilibrium potential or capacitary potential. We will extend it to $`F`$ by setting $`P_F(x)=1`$ for all $`xF`$.
It follows from the maximum principle that $`0P_F1`$ everywhere in $`\text{}^n`$.
In case when $`F`$ is the closure of an open subset with a smooth boundary, $`u=P_F`$ is the unique minimizer for the Dirichlet integral in (2.1). In particular,
$$|P_F|^2𝑑x=\mathrm{cap}(F).$$
where the integration is taken over $`\text{}^n`$ (or $`\text{}^nF`$).
The capacity of the ball $`\overline{B}_r`$ is easily calculated and is given by
(2.4)
$$\mathrm{cap}(\overline{B}_r)=(n2)\omega _nr^{n2}.$$
## 3 Lower bound
In this section we will establish the lower bound for $`\lambda (\mathrm{\Omega })`$ from Theorem 1.1 which is an easier part of this theorem. The key part of the lower bound proof is presented in the following lemma, which was first proved in (see also , where it is present as a particular case of a much more general Theorem 10.1.2, part 1), though without an explicit constant, which we provide to specify explicit constants in Theorem 1.1.
###### Lemma 3.1
The following inequality holds for every complex-valued function $`u\mathrm{Lip}(\overline{B}_r)`$ which vanishes on a compact set $`F\overline{B}_r`$ (but is not identically $`0`$ on $`\overline{B}_r`$):
(3.1)
$$\mathrm{cap}(F)\frac{C_n_{B_r}|u(x)|^2𝑑x}{r^n_{B_r}|u(x)|^2𝑑x},$$
where
(3.2)
$$C_n=4\omega _n\left(1\frac{2}{n^2}\right)$$
Beginning of Proof. A. Clearly, it is sufficient to consider the ball $`B_r`$ centered at $`0`$, and real-valued functions $`u\mathrm{Lip}(\overline{B}_r)`$. By scaling we see that it suffices to consider the case $`r=1`$. (The corresponding estimate for an arbitrary $`r>0`$ follows from the one with $`r=1`$ with the same constant $`C_n`$.) So we need to prove the estimate
(3.3)
$$_{B_1}|u|^2𝑑x\frac{C_n}{\mathrm{cap}(F)}_{B_1}|u|^2𝑑x,$$
where $`F`$ is a compact subset of $`\overline{B}_1`$, $`u\mathrm{Lip}(\overline{B}_1)`$ and $`u|_F=0`$.
To be able to use (2.1), consider the following function $`U\mathrm{Lip}(\text{}^n)`$:
$$U(x)=\{\begin{array}{cc}1|u(x)|,\hfill & \text{if }|x|1\text{,}\hfill \\ |x|^{2n}(1|u(|x|^2x)|),\hfill & \text{if }|x|1,\hfill \end{array}$$
i.e. $`U`$ extends $`1|u|`$ to $`\{x:|x|1\}`$ as the Kelvin transform of $`1|u|`$. Clearly, $`U|_F=1`$, $`|U|=|u|`$ almost everywhere in $`B_1`$, $`U(x)=O(|x|^{2n})`$ and $`|U(x)|=O(|x|^{1n})`$ as $`|x|\mathrm{}`$. It follows that $`U`$ can serve as a test function in (2.1), i.e.
(3.4)
$$\mathrm{cap}(F)_\text{}^n|U|^2𝑑x.$$
Using the harmonicity of $`|x|^{2n}`$ and the Green-Stokes formula, we obtain by a straightforward calculation
(3.5)
$$_\text{}^n|U|^2𝑑x=2_{B_1}|u|^2𝑑x+(n2)_{B_1}(1|u(\omega )|)^2𝑑\omega ,$$
where $`B_1=\{\omega \text{}^n,|\omega |=1\}`$ is the boundary of $`B_1`$ (the unit sphere in $`\text{}^n`$), $`d\omega `$ means the standard volume element on $`B_1`$.
B. For a function $`v`$ on $`B_1`$ define its average
$$\overline{v}=_{B_1}vd\omega =\frac{1}{\omega _n}_{B_1}v𝑑\omega .$$
To continue the proof of Lemma 3.1, we will need the following elementary lemma.
###### Lemma 3.2
For any $`v\mathrm{Lip}(B_1)`$,
(3.6)
$$_{B_1}|v\overline{v}|^2𝑑\omega _{B_1}|v|^2𝑑x.$$
Proof of Lemma 3.2. It suffices to prove it for real-valued functions $`v`$. Let us expand $`v`$ in spherical functions. Let
$$\{Y_{k,l}|l=0,1,\mathrm{},n_k,k=0,1,\mathrm{}\}$$
be an orthonormal basis in $`L^2(B_1)`$ which consists of eigenfunctions of the (negative) Laplace-Beltrami operator $`\mathrm{\Delta }_\omega `$ on $`B_1`$, so that the eigenfunctions $`Y_{k,l}=Y_{k,l}(\omega )`$ with a fixed $`k`$ have the same eigenvalue $`k(k+n2)`$ (which has multiplicity $`n_k+1`$). Note that the zero eigenvalue (corresponding to $`k=0`$) has multiplicity 1 and $`Y_{0,0}=const=\omega _n^{1/2}`$ for the corresponding eigenfunction.
Writing $`x=r\omega `$, where $`r=|x|`$, $`\omega =x/|x|`$, we can present $`v`$ in the form
(3.7)
$$v(x)=v(r,\omega )=\underset{k,l}{}v_{k,l}(r)Y_{k,l}(\omega ).$$
Then
(3.8)
$$_{B_1}|v(x)|^2𝑑x=\underset{k,l}{}_0^1|v_{k,l}(r)|^2r^{n1}𝑑r,$$
and
(3.9)
$$_{B_1}|v(\omega )|^2𝑑\omega =\underset{k,l}{}|v_{k,l}(1)|^2.$$
It follows that
(3.10)
$$_{B_1}|v(\omega )\overline{v}|^2𝑑\omega =\underset{\{k,l:k1\}}{}|v_{k,l}(1)|^2.$$
Taking into account that
$$|v|^2=\left|\frac{v}{r}\right|^2+r^2|_\omega v|^2,$$
where $`_\omega `$ means the gradient along the unit sphere with variable $`\omega `$ and fixed $`r`$, we also get
(3.11)
$$_{B_1}|v|^2𝑑x=\underset{k,l}{}_0^1\left(|v_{k,l}^{}(r)|^2+\frac{k(k+n2)}{r^2}|v_{k,l}(r)|^2\right)r^{n1}𝑑r.$$
Comparing (3.10) and (3.11), and taking into account that $`k(k+n2)`$ increases with $`k`$, we see that it suffices to establish that the inequality
$$|g(1)|^2_0^1\left(|g^{}(r)|^2+\frac{n1}{r^2}|g(r)|^2\right)r^{n1}𝑑r,$$
holds for any real-valued function $`g\mathrm{Lip}([0,1]))`$. To this end write
$`g(1)^2`$ $`={\displaystyle _0^1}(r^{n2}g^2)^{}𝑑r={\displaystyle _0^1}[2r^{n2}g^{}g+(n2)r^{n3}g^2]𝑑r`$
$`{\displaystyle _0^1}[r^{n1}g^2+(n1)r^{n3}g^2]𝑑r={\displaystyle _0^1}\left(g^2+{\displaystyle \frac{n1}{r^2}}g^2\right)r^{n1}𝑑r,`$
which proves Lemma 3.2. $`\mathrm{}`$
C. Proof of Lemma 3.1 (continuation). Let us normalize $`u`$ by requiring $`\overline{|u|}=1`$, i.e. average of $`|u|`$ over $`B_1`$ equals 1. This can be done if $`u0`$ on $`B_1`$. Then by Lemma 3.2
$$_{B_1}(1|u|)^2𝑑\omega _{B_1}|u|^2𝑑x.$$
Combining this with (3.4) and (3.5), we obtain
$$\mathrm{cap}(F)n_{B_1}|u|^2𝑑x.$$
Removing the restriction $`\overline{|u|}=1`$, we can conclude that for any $`u\mathrm{Lip}(B_1)`$
(3.12)
$$\left(_{B_1}|u|d\omega \right)^2\frac{n}{\mathrm{cap}(F)}_{B_1}|u|^2𝑑x.$$
(This obviously also holds in case when $`u0`$ on $`B_1`$.)
Note that for any real-valued function $`v\mathrm{Lip}(B_1)`$
$$_{B_1}|v\overline{v}|^2d\omega =_{B_1}|v|^2d\omega \overline{v}^2,$$
hence, using (3.6), we get
$$_{B_1}|v|^2d\omega =\overline{v}^2+_{B_1}|v\overline{v}|^2d\omega \overline{v}^2+\frac{1}{\omega _n}_{B_1}|v|^2𝑑x.$$
Applying this to $`v=|u|`$ and using (3.12), we obtain
(3.13)
$$_{B_1}|u|^2𝑑\omega \left(1+\frac{n\omega _n}{\mathrm{cap}(F)}\right)_{B_1}|u|^2𝑑x.$$
D. Note that out goal is an estimate which is similar to (3.13) but with the integral over $`B_1`$ in the left hand side replaced by the integral over $`B_1`$. To this end we will again use the expansion (3.7) of $`v=|u|`$ over spherical functions, and the identities (3.8), (3.9) and (3.11). Let us take a real-valued function $`g\mathrm{Lip}([0,1])`$ and denote
$$Q=_0^1g^2(r)r^{n1}𝑑r.$$
Integrating by parts, we obtain
$$Q=\frac{2}{n}_0^1gg^{}r^n𝑑r+\frac{1}{n}g^2(1).$$
Using an elementary inequality $`2ab\epsilon a^2+\epsilon ^1b^2`$, where $`a,b\text{}`$, $`\epsilon >0`$, and taking into account that $`r1`$, we obtain
$`Q`$ $`{\displaystyle \frac{1}{n}}{\displaystyle _0^1}\left(\epsilon g^2(r)+{\displaystyle \frac{1}{\epsilon }}g^2(r)\right)r^{n1}𝑑r+{\displaystyle \frac{1}{n}}g^2(1)`$
$`={\displaystyle \frac{\epsilon }{n}}Q+{\displaystyle \frac{1}{n\epsilon }}{\displaystyle _0^1}g^2(r)r^{n1}𝑑r+{\displaystyle \frac{1}{n}}g^2(1),`$
hence for any $`\epsilon (0,n)`$
$$Q\frac{1}{(n\epsilon )\epsilon }_0^1g^2(r)r^{n1}𝑑r+\frac{1}{n\epsilon }g^2(1).$$
Taking $`\epsilon =n/2`$, we obtain
(3.14)
$$Q\frac{4}{n^2}_0^1g^2(r)r^{n1}𝑑r+\frac{2}{n}g^2(1).$$
Now we can argue as in the proof of Lemma 3.2, expanding $`v=|u|`$ over spherical harmonics $`Y_{k,l}`$. Then the desired inequality follows from the inequalities for the coefficients $`v_{k,l}=v_{k,l}(r)`$, with the strongest one corresponding to the case $`k=0`$ (unlike $`k=1`$ in Lemma 3.2). Then using the inequality (3.14) for $`g=v_{0,0}`$ we obtain
(3.15)
$$_{B_1}|u|^2𝑑x\frac{4}{n^2}_{B_1}|u|^2𝑑x+\frac{2}{n}_{B_1}|u|^2𝑑\omega .$$
Using (3.13), we deduce from (3.15):
(3.16)
$$_{B_1}|u|^2𝑑x\left[\frac{4}{n^2}+\frac{2}{n}\left(1+\frac{n\omega _n}{\mathrm{cap}(F)}\right)\right]_{B_1}|u|^2𝑑x.$$
Taking into account the inequality
$$\mathrm{cap}(F)\mathrm{cap}(\overline{B}_1)=(n2)\omega _n,$$
we can estimate the constant in front of the integral in the right hand side of (3.16) as follows:
$$\frac{4}{n^2}+\frac{2}{n}\left(1+\frac{n\omega _n}{\mathrm{cap}(F)}\right)\frac{4\omega _n}{\mathrm{cap}(F)}\left(1\frac{2}{n^2}\right),$$
which ends the proof of Lemma 3.1. $`\mathrm{}`$
The lower bound in (1.3) is given by
###### Lemma 3.3
There exists $`c=c(\gamma ,n)>0`$ such that for all open sets $`\mathrm{\Omega }\text{}^n`$
(3.17)
$$\lambda (\mathrm{\Omega })cr_{\mathrm{\Omega },\gamma }^2.$$
Proof. Let us fix $`\gamma (0,1)`$ and choose any $`r>r_{\mathrm{\Omega },\gamma }`$. Then any ball $`\overline{B}_r`$ has a non-negligible intersection with $`\text{}^n\mathrm{\Omega }`$, i.e.
$$\mathrm{cap}(\overline{B}_r\mathrm{\Omega })\gamma \mathrm{cap}(\overline{B}_r).$$
Since any $`uC_0^{\mathrm{}}(\mathrm{\Omega })`$ vanishes on $`\overline{B}_r\mathrm{\Omega }`$, it follows from Lemma 3.1 that for any such $`u`$
$$_{\overline{B}_r}|u|^2𝑑x\frac{C_n}{r^n\mathrm{cap}(\overline{B}_r\mathrm{\Omega })}_{\overline{B}_r}|u|^2𝑑x\frac{C_n}{r^n\gamma \mathrm{cap}(\overline{B}_r)}_{\overline{B}_r}|u|^2𝑑x.$$
Taking into account that $`\mathrm{cap}(\overline{B}_r)=\mathrm{cap}(\overline{B}_1)r^{n2}`$, we obtain
$$_{\overline{B}_r}|u|^2𝑑x\frac{C_nr^2}{\gamma \mathrm{cap}(\overline{B}_1)}_{\overline{B}_r}|u|^2𝑑x.$$
Now let us choose a covering of $`\text{}^n`$ by balls $`\overline{B}_r=\overline{B}_r^{(k)}`$, $`k=1,2,\mathrm{},`$ so that the multiplicity of this covering is at most $`N=N(n)`$. For example, we can make
(3.18)
$$N(n)n\mathrm{log}n+n\mathrm{log}(\mathrm{log}n)+5n,n2,$$
which holds also for the smallest multiplicity of coverings of $`\text{}^n`$ by translations of any convex body (see Theorem 3.2 in ).
Then summing up the estimates above over all balls in this covering, we see that
$`{\displaystyle _\text{}^n}|u|^2𝑑x{\displaystyle \underset{k}{}}{\displaystyle _{\overline{B}_r^{(k)}}}|u|^2𝑑x{\displaystyle \frac{C_nr^2}{\gamma \mathrm{cap}(\overline{B}_1)}}{\displaystyle \underset{k}{}}{\displaystyle _{\overline{B}_r^{(k)}}}|u|^2𝑑x`$
$`{\displaystyle \frac{C_nNr^2}{\gamma \mathrm{cap}(\overline{B}_1)}}{\displaystyle _\text{}^n}|u|^2𝑑x.`$
Recalling (1.1), we see that
$$\lambda (\mathrm{\Omega })cr^2$$
with
(3.19)
$$c=c(\gamma ,n)=\frac{\gamma \mathrm{cap}(\overline{B}_1)}{C_nN}=\frac{\gamma n^2(n2)}{4(n^22)N}.$$
Taking limit as $`rr_{\mathrm{\Omega },\gamma }`$, we obtain (3.17) with the same $`c`$. $`\mathrm{}`$
## 4 Upper bound
### 4.1
According to (1.1), to get an upper bound for $`\lambda (\mathrm{\Omega })`$ it is enough to take any test function $`uC_0^{\mathrm{}}(\mathrm{\Omega })`$ and write
(4.1)
$$\lambda (\mathrm{\Omega })\frac{_\mathrm{\Omega }|u|^2𝑑x}{_\mathrm{\Omega }|u|^2𝑑x}.$$
For simplicity of notations we will write $`\lambda `$ instead of $`\lambda (\mathrm{\Omega })`$ everywhere in this section. The inequality (4.1) can be rewritten as follows:
(4.2)
$$_\mathrm{\Omega }|u|^2𝑑x\lambda ^1_\mathrm{\Omega }|u|^2𝑑x.$$
By approximation, it suffices to take $`u\mathrm{Lip}_c(\mathrm{\Omega })`$ or even $`uH_0^1(\mathrm{\Omega })`$, where $`H_0^1(\mathrm{\Omega })`$ is the closure of $`C_0^{\mathrm{}}(\mathrm{\Omega })`$ in the standard Sobolev space $`H^1(\mathrm{\Omega })`$ (which consists of all $`uL^2(\mathrm{\Omega })`$ with the distributional derivatives $`u/x_jL^2(\mathrm{\Omega })`$, $`j=1,\mathrm{},n`$).
In particular, choosing a ball $`B_r`$, we can take
(4.3)
$$u\mathrm{Lip}_c(\mathrm{\Omega }B_r)=\mathrm{Lip}_c(\mathrm{\Omega })\mathrm{Lip}_c(B_r).$$
Let us take a compact set $`F\overline{B}_{3r/2}`$, such that $`F`$ is the closure of an open set with a smooth boundary. (In this section we will call such sets regular subsets of $`\overline{B}_{3r/2}`$.) Denote by $`P_F`$ its equilibrium potential (see Sect. 2). Regularity of $`F`$ implies that $`P_F\mathrm{Lip}(\text{}^n)`$. By definition $`P_F=1`$ on $`F`$, so $`1P_F=0`$ on $`F`$. Let us also assume that
$$\mathrm{Int}F\overline{B}_r\mathrm{\Omega },$$
where $`\mathrm{Int}F`$ means the set of all interior points of $`F`$ (so $`\mathrm{Int}F`$ is an open subset in $`\text{}^n`$). Then $`1P_F=0`$ in a neighborhood of $`\overline{B}_r\mathrm{\Omega }`$. Therefore, multiplying $`1P_F`$ by a cut-off function $`\eta C_0^{\mathrm{}}(B_r)`$, we will get a function $`u=\eta (1P_F)`$, satisfying the requirenment (4.3), hence fit to be a test function in (4.1).
In the future we will also assume that the cut-off function $`\eta C_0^{\mathrm{}}(B_r)`$ has the following properties:
$$0\eta 1\mathrm{on}B_r,\eta =1\mathrm{on}B_{(1\kappa )r},|\eta |2(\kappa r)^1\mathrm{on}B_r,$$
where $`0<\kappa <1`$ and the balls $`B_r`$ and $`B_{(1\kappa )r}`$ are supposed to have the same center. Using integration by parts and the equation $`\mathrm{\Delta }P_F=0`$ on $`B_rF`$, we obtain for the test function $`u=\eta (1P_F)`$:
$`{\displaystyle _{B_r}}|u|^2𝑑x`$ $`={\displaystyle _{B_r}}\left(|\eta |^2(1P_F)^2(\eta ^2)(1P_F)P_F+\eta ^2|P_F|^2\right)𝑑x`$
$`={\displaystyle _{B_r}}|\eta |^2(1P_F)^2𝑑x4(\kappa r)^2{\displaystyle _{B_r}}(1P_F)^2𝑑x.`$
Therefore, from (4.2) we obtain
$$_{B_r}|u|^2𝑑x\lambda ^14(\kappa r)^2_{B_r}(1P_F)^2𝑑x.$$
Since $`0P_F1`$, the last integral in the right hand side is estimated by
$$\mathrm{mes}(B_r)=n^1\omega _nr^n.$$
where $`\mathrm{mes}`$ means the usual Lebesgue measure on $`\text{}^n`$. Therefore,
$$_{B_r}|u|^2𝑑x4n^1\omega _n\lambda ^1\kappa ^2r^{n2}.$$
Restricting the integral in the left hand side to $`B_{(1\kappa )r}`$, we obtain
(4.4)
$$_{B_{(1\kappa )r}}(1P_F)^2𝑑x4n^1\omega _n\lambda ^1\kappa ^2r^{n2}.$$
### 4.2
Now we need to provide an appropriate lower bound for the left hand side of (4.4). To this end we first restrict the integration to the spherical layer
$$S_{r_1,r_2}=B_{r_2}B_{r_1},$$
where $`0<r_1<r_2<r`$. In the future we will take
(4.5)
$$r_1=(12\kappa )r,r_2=(1\kappa )r,$$
where $`0<\kappa <1/2`$, though it is convenient to write some formulas in a bigger generality. Let us denote the volume of the layer $`S_{r_1,r_2}`$ by $`|S_{r_1,r_2}|`$, i.e.
$$|S_{r_1,r_2}|=\mathrm{mes}S_{r_1,r_2}=n^1\omega _n(r_2^nr_1^n).$$
We will also need the notation
$$_{S_{r_1,r_2}}f(x)dx=\frac{1}{|S_{r_1,r_2}|}_{S_{r_1,r_2}}f(x)𝑑x$$
for the average of a positive function $`f`$ over $`S_{r_1,r_2}`$. In particular, restricting the integration in (4.4) to $`S_{r_1,r_2}`$ (with $`r_1,r_2`$ as in (4.5)) and dividing by $`|S_{r_1,r_2}|`$, we obtain
$$_{S_{r_1,r_2}}(1P_F)^2dx\frac{4\lambda ^1\kappa ^2r^{n2}}{r_2^nr_1^n}.$$
Hence, by the Cauchy-Schwarz inequality,
(4.6)
$$\left[1_{S_{r_1,r_2}}P_Fdx\right]^2=\left[_{S_{r_1,r_2}}(1P_F)dx\right]^2\frac{4\lambda ^1\kappa ^2r^{n2}}{r_2^nr_1^n}.$$
To simplify the right hand side, let us estimate $`(r_2^nr_1^n)^1`$ from above. Applying the Bernoulli inequality, we see that
$`r_2^nr_1^n=(r_2r_1)(r_2^{n1}+r_2^{n2}r_1+\mathrm{}+r_1^{n1})`$
$`n\kappa rr_1^{n1}=n\kappa r^n(12\kappa )^{n1}n\kappa r^n[12(n1)\kappa ].`$
Now note that
$$\frac{1}{12(n1)\kappa }1+4(n1)\kappa ,$$
provided
(4.7)
$$0<\kappa \frac{1}{4(n1)}.$$
Under this condition it follows that
(4.8)
$$\frac{1}{r_2^nr_1^n}n^1\kappa ^1r^n\left[1+4(n1)\kappa \right],$$
and (4.6) takes the form
(4.9)
$$\left[1_{S_{r_1,r_2}}P_Fdx\right]^24n^1\kappa ^3\left[1+4(n1)\kappa \right]\lambda ^1r^2.$$
### 4.3
For simplicity of notations and without loss of generality we may assume that the ball $`B_r`$ is centered at $`0\text{}^n`$ (and so are smaller balls and spherical layers).
To provide a lower bound for the left hand side of (4.9), we will give an upper bound for the average of $`P_F`$. According to the definition of $`P_F`$ and notations from Section 2, we can write
(4.10) $`{\displaystyle }_{S_{r_1,r_2}}P_Fdx`$ $`={\displaystyle }_{S_{r_1,r_2}}\left({\displaystyle _F}(xy)𝑑\mu _F(y)\right)dx`$
$`={\displaystyle _F}\left({\displaystyle }_{S_{r_1,r_2}}(xy)dx\right)𝑑\mu _F(y).`$
The inside integral in the right hand side can be explicitly calculated (due to Newton) as the potential of a uniformly charged spherical layer with total charge 1. The result of this calculation is $`|S_{r_1,r_2}|^1V_{r_1,r_2}(y)`$, where
(4.11) $`V_{r_1,r_2}(y)=\{\begin{array}{cc}\frac{r_2^2r_1^2}{2(n2)},\hfill & \text{if }|y|r_1,\hfill \\ \frac{|y|^2}{2n}+\frac{r_2^2}{2(n2)}\frac{r_1^n}{n(n2)|y|^{n2}},\hfill & \text{if }r_1|y|r_2,\hfill \\ \frac{r_2^nr_1^n}{n(n2)|y|^{n2}},\hfill & \text{if }|y|r_2.\hfill \end{array}`$
The function $`yV_{r_1,r_2}(y)`$ belongs to $`C^1(\text{}^n)`$ and is spherically symmetric; it tends to $`0`$ as $`|y|\mathrm{}`$; it is harmonic in $`\text{}^nS_{r_1,r_2}`$ and satisfies the equation $`\mathrm{\Delta }V_{r_1,r_2}=1`$ in $`S_{r_1,r_2}`$. These properties uniquely define the function $`V_{r_1,r_2}`$. Differentiating it with respect to $`|y|`$, we easily see that it is decreasing with respect to $`|y|`$, hence its maximum is at $`y=0`$ (hence given by the first row in (4.11)). So we obtain, using (4.8):
$`{\displaystyle }_{S_{r_1,r_2}}(xy)dx|S_{r_1,r_2}|^1V_{r_1,r_2}(0)={\displaystyle \frac{n(r_2^2r_1^2)}{2(n2)\omega _n(r_2^nr_1^n)}}`$
$`={\displaystyle \frac{n\kappa r(r_1+r_2)}{2(n2)\omega _n(r_2^nr_1^n)}}{\displaystyle \frac{nr^2\kappa (1\kappa )}{(n2)\omega _n(r_2^nr_1^n)}}`$
$`{\displaystyle \frac{(1\kappa )[1+4(n1)\kappa ]}{(n2)\omega _nr^{n2}}}{\displaystyle \frac{1+(4n5)\kappa }{(n2)\omega _nr^{n2}}}`$
Finally, using (2.4), we obtain
(4.12)
$$_{S_{r_1,r_2}}(xy)dx\frac{1+(4n5)\kappa }{\mathrm{cap}(\overline{B}_r)}.$$
provided $`r_1,r_2`$ choosen as in (4.5) and (4.7) is satisfied.
### 4.4
Using (4.12) in (4.10) and taking into account (2.3), we obtain
(4.13) $`{\displaystyle }_{S_{r_1,r_2}}P_F(x)dx{\displaystyle \frac{1+(4n5)\kappa }{\mathrm{cap}(\overline{B}_r)}}{\displaystyle _F}𝑑\mu _F(y)`$
$`=[(1+(4n5)\kappa ]{\displaystyle \frac{\mathrm{cap}(F)}{\mathrm{cap}(\overline{B}_r)}}[(1+(4n5)\kappa ]\gamma ,`$
provided $`F`$ is $`\gamma `$-negligible. (i.e. satisfies (1.2)). Note that we do not assume that $`F\overline{B}_r`$ but do assume that $`0<\gamma <1`$. In this case, taking into account (4.7), we can take
(4.14)
$$\kappa =\mathrm{min}\{\frac{1}{4(n1)},\frac{1\gamma }{2(4n5)\gamma }\},$$
so that (4.7) is satisfied, and, besides,
$$[(1+(4n5)\kappa ]\gamma \frac{1+\gamma }{2}=1\frac{1\gamma }{2},$$
so that (4.13) becomes
$$_{S_{r_1,r_2}}P_F(x)dx1\frac{1\gamma }{2}.$$
Taking this into account in (4.9) and using (4.7), we obtain
$$\frac{(1\gamma )^2}{4}4n^1\kappa ^3[1+4(n1)\kappa ]\lambda ^1r^28n^1\kappa ^3\lambda ^1r^2,$$
hence
(4.15)
$$\lambda 32(1\gamma )^2\kappa ^3r^2.$$
### 4.5
We are now ready for
Proof of Theorem 1.1.
The lower bound for $`\lambda `$ was established in Lemma 3.3.
We proved the estimate (4.15) above under the condition that there exist $`\gamma (0,1)`$, a ball $`B_r`$ and a regular compact set $`F\overline{B}_{3r/2}`$ (here the balls $`B_r`$ and $`B_{3r/2}`$ have the same center), such that $`F`$ is $`\gamma `$-negligible and its interior includes $`\overline{B}_r\mathrm{\Omega }`$. (The estimate then holds with $`\kappa =\kappa (\gamma ,n)`$ given by (4.14).) It follows in particular that $`\overline{B}_r\mathrm{\Omega }`$ is $`\gamma `$-negligible.
Conversely, if $`\overline{B}_r\mathrm{\Omega }`$ is $`\gamma `$-negligible, then we can approximate it by regular compact sets $`F_k`$, $`k=1,2,\mathrm{}`$, such that $`\mathrm{Int}F_k\overline{B}_r\mathrm{\Omega }`$, $`\mathrm{Int}F_kF_{k+1}`$, and $`\overline{B}_r\mathrm{\Omega }`$ is the intersection of all $`F_k`$’s. In this case
$$\underset{k\mathrm{}}{lim}\mathrm{cap}(F_k)=\mathrm{cap}(\overline{B}_r\mathrm{\Omega }),$$
due to the continuity property of the capacity. (See e.g. , Sect. 2.2.1.) In this case, for any $`\epsilon >0`$ the sets $`F_k`$ will be $`(\gamma +\epsilon )`$-negligible for sufficiently large $`k`$. It follows that the estimate (4.15) will hold if we only know that there exists a ball $`B_r`$ such that $`\overline{B}_r\mathrm{\Omega }`$ is $`\gamma `$-negligible. Then the estimate still holds if we replace $`r`$ by the least upper bound of the radii of such balls which is exactly the interior capacitary radius $`r_{\mathrm{\Omega },\gamma }`$. This proves the upper bound in (1.3) with
(4.16)
$$C(\gamma ,n)=32(1\gamma )^2\kappa ^3,$$
where $`\kappa `$ is defined by (4.14). $`\mathrm{}`$
## 5 Further remarks
### 5.1 Measure instead of capacity
E. Lieb used geometric arguments to establish a lower bound for $`\lambda (\mathrm{\Omega })`$ which is similar to (3.17) but with capacity replaced by the Lebesgue measure. Such lower bounds can be also deduced from Theorem 1.1 if we use isoperimetric inequalities between the capacity and Lebesgue measure:
(5.1)
$$\mathrm{mes}FA_n(\mathrm{cap}(F))^{n/(n2)},$$
with the equality for balls (see e.g. or Sect. 2.2.3, 2.2.4 in ), so
$$A_n=(\mathrm{mes}B_1)\left[\mathrm{cap}(B_1)\right]^{n/(n2)}=n^1(n2)^{n/(n2)}\omega _n^{2/(n2)}.$$
Namely, let us denote for any $`\alpha (0,1)`$
$$r_{\mathrm{\Omega },\alpha }^{(mes)}=sup\{r|B_r,\mathrm{mes}(B_r\mathrm{\Omega })\alpha \mathrm{mes}B_r\}.$$
Then (5.1) implies that
$$r_{\mathrm{\Omega },\alpha }^{(mes)}r_{\mathrm{\Omega },\gamma }\text{provided}\alpha =\gamma ^{n/(n2)}.$$
Therefore, we obtain
###### Proposition 5.1
For every $`\alpha (0,1)`$
$$\lambda (\mathrm{\Omega })c(\gamma ,n)\left(r_{\mathrm{\Omega },\alpha }^{(mes)}\right)^2,\text{where}\gamma =\alpha ^{(n2)/n}.$$
Here $`c(\gamma ,n)`$ is given by (3.19).
This is exactly Lieb’s inequality (1.2) in , though with a different constant.
There are numerous results which give lower bounds for $`\lambda (\mathrm{\Omega })`$. We will mention only a few. The famous Faber-Krahn inequality () gives a lower bound of $`\lambda (\mathrm{\Omega })`$ in terms of the area of $`\mathrm{\Omega }\text{}^2`$. Under miscellaneous topological and geometric restrictions on $`\mathrm{\Omega }`$ the interior radius was shown to provide a lower bound (hence a two-sided estimate) for $`\lambda (\mathrm{\Omega })`$ in case $`n=2`$ by Hayman , Osserman , Taylor , Croke , Bañuelos and Carroll , and also in case $`n3`$ (, ).
Let $`\mathrm{cap}_\mathrm{\Omega }(F)`$ denote the capacity of a compact set $`F\mathrm{\Omega }`$ with respect to an open set $`\mathrm{\Omega }\text{}^n`$. It is defined similarly to $`\mathrm{cap}(F)`$ in (2.1), except the allowed test functions $`u`$ should be supported in $`\mathrm{\Omega }`$. The following 2-sided estimate for $`\lambda (\mathrm{\Omega })`$ was established in :
(5.2)
$$\frac{1}{4}\underset{F}{inf}\frac{\mathrm{cap}_\mathrm{\Omega }(F)}{\mathrm{mes}F}\lambda (\mathrm{\Omega })\underset{F}{inf}\frac{\mathrm{cap}_\mathrm{\Omega }(F)}{\mathrm{mes}F},$$
where the infimum is taken over all compact sets $`F\mathrm{\Omega }`$. The constant 1/4 in the lower bound is precise. Both inequalities in (5.2) hold on Riemannian manifolds as well.
The lower bound (the first inequality) in (5.2) implies the Cheeger inequality (see Sect. 3 in for this implication), which also provides a geometric lower bound for $`\lambda (\mathrm{\Omega })`$ on manifolds. (See also Grigor’yan for a review and related results.)
### 5.2 Bounds for essential spectrum
Let $`\lambda _{\mathrm{}}(\mathrm{\Omega })`$ denote the bottom of the essential spectrum of $`\mathrm{\Delta }`$ with the Dirichlet boundary conditions in $`\mathrm{\Omega }`$. Then Persson’s arguments give
$$\lambda _{\mathrm{}}(\mathrm{\Omega })=\underset{R+\mathrm{}}{lim}\lambda (\mathrm{\Omega }\overline{B}_R(0)),$$
where $`B_R(0)`$ is the ball with the radius $`R`$ and the center at the origin. Applying two-sided estimates from Theorem 1.1 to $`\lambda (\mathrm{\Omega }\overline{B}_R(0))`$, we obtain
###### Theorem 5.2
For any $`\gamma (0,1)`$ and any open set $`\mathrm{\Omega }\text{}^n`$,
$$cr_{\mathrm{\Omega },\gamma ,\mathrm{}}^2\lambda _{\mathrm{}}(\mathrm{\Omega })Cr_{\mathrm{\Omega },\gamma ,\mathrm{}}^2,$$
where
$$r_{\mathrm{\Omega },\gamma ,\mathrm{}}=\underset{R\mathrm{}}{lim}r_{\mathrm{\Omega }\overline{B}_R(0),\gamma },$$
and the constants $`c=c(\gamma ,n)`$, $`C=C(\gamma ,n)`$ are the same as in Theorem 1.1.
For small $`\gamma `$ this theorem is due to Maz’ya and Otelbaev (see and also Theorem 12.3.1 in ).
Theorem 5.2 implies that for any $`\gamma (0,1)`$ the condition $`r_{\mathrm{\Omega },\gamma ,\mathrm{}}=0`$ is necessary and sufficient for the discreteness of spectrum of the operator $`\mathrm{\Delta }`$ with the Dirichlet boundary conditions in $`\mathrm{\Omega }`$. (This is also a particular case of the main results of ).
### 5.3 Bounds for spectra of Schrödinger operators
Theorem 1.1, Proposition 5.1 and Theorem 5.2 can be extended to Schrödinger operators with positive potentials (which are even allowed to be positive measures, which are absolutely continuous with respect to the Wiener capacity). For small $`\gamma `$ these results can be found in Chapters 10 – 12 of with appropriate references.
For simplicity of formulation we will consider operators $`H_V=\mathrm{\Delta }+V`$, $`V0`$, where $`V`$ is locally integrable. Then 2-sided estimates of the type (1.3) can be obtained for the bottom of the spectrum (and essential spectrum) of $`H_V`$, if $`r_{\mathrm{\Omega },\gamma }`$ is replaced by the quantity
$$r_{V,\gamma }=sup\left\{r\right|B_r,\text{such that}r^{n2}\underset{F}{inf}_{B_rF}V𝑑x\},$$
where the infimum is taken over all negligible subsets in $`F\overline{B}_r`$, i.e. sets satisfying (1.2).
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# Quantum error correction of systematic errors using a quantum search framework
## I Introduction
Quantum error correction is perhaps the biggest hurdle in building a quantum computer. Imperfect control operations are one of several sources of error.While error-correction schemes designed to correct for general errors Preskill97 no doubt also correct control errors, error-correction or error-avoidance schemes tuned to the dominant physical error model are more efficient and practical. Specialized error-correction schemes can also tolerate higher noise rates.
This paper specializes to systematic control errors of the following form. When we try to apply the single-qubit pulse $`U=\mathrm{exp}(i\theta \widehat{n}\stackrel{}{\sigma })`$, a $`2\theta `$ rotation about axis $`\widehat{n}`$, we in fact apply $`\overline{U}=UV`$. The error is systematic in the sense that it is invertible; attempting to apply $`U^{}`$ in fact applies $`\overline{U}^{}`$. The form of the error $`V`$ is of course restricted. Previous authors McHughTwamley04 ; BrownHarrowChuang04 have considered the case of linear over-rotational errors: $`V=\mathrm{exp}(iϵ\theta \widehat{n}\stackrel{}{\sigma })`$ where $`ϵ`$ is fixed and small, but unknown. Here, we consider the case of general over-rotational errors, $`V=\mathrm{exp}[iϵ(|\theta |,\widehat{n})\widehat{n}\stackrel{}{\sigma }]`$. The amount of over-rotation, $`ϵ(|\theta |,\widehat{n})`$, can now depend arbitrarily on the rotation angle $`2\theta `$ and also the axis of rotation $`\widehat{n}`$.
Our error-correction method is a new composite pulse sequence, inspired by the generalization of quantum search known as amplitude amplification BrassardHoyerTapp98 . In this algorithm, the amplitude produced in a particular target subspace by applying some unitary $`\overline{U}`$ to a source state is amplified by successively repeating $`R_0(\pi )\overline{U}^{}R_t(\pi )\overline{U}`$. Here $`R_0(\pi )`$ and $`R_t(\pi )`$ are selective reflections about the source and target, respectively. In standard quantum search, the source is $`|0^n`$, $`\overline{U}=H^n`$ transverse Hadamard, and the target is a bit string $`|x`$. In the subspace spanned by $`\overline{U}|0^n`$ and $`|x`$, the state vector steadily rotates toward $`|x`$. Eventually, it rotates past the target.
What happens in quantum search if we don’t merely reflect about the source and target, but instead add a phase other than $`\pi `$? It is well known that any phase bounded away from 0 works to give a square-root speedup (with different constants); for example, this fact is used in a stronger form in Ambainis’s element distinctness algorithm Ambainis03 . One of us (LG) noticed that concatenating the basic sequence
$$\overline{U}R_0(\pi /3)\overline{U}^{}R_t(\pi /3)\overline{U}$$
results in the state converging to the target subspace and not overshooting it, when viewed at times $`3^k`$, $`k𝐍`$ Grover05 . Figure 1 gives some geometrical intuition.
For the present problem of systematic control errors, there is no source or target – we desire a “fully compensating” pulse sequence accurate on an arbitrary input – but a similar calculation still applies. We need merely choose a source arbitrarily, say the $`X`$ $`+1`$ eigenstate $`|+`$, and set the target accordingly (to $`U|+`$). Assume $`\widehat{n}=(0,0,1)`$. Let $`R_0=\mathrm{exp}(i\frac{\pi }{6}X)`$ a $`\pi /3`$ rotation about the $`X`$ axis, and $`R_t=UR_0U^{}`$. When we apply the sequence of noisy pulses
$$\overline{R}_t^{}\overline{U}\overline{R}_0\overline{U}^{}\overline{R}_t\overline{U}\overline{R}_0^{},$$
the different systematic errors in both $`\overline{U}`$ and the noisy $`\pi /3`$ rotations largely cancel out, leaving behind a higher-order error. (The extra pulses at either end adjust for the phase difference between $`|+`$ and its orthogonal complement $`|`$ which would otherwise be introduced.) This correction sequence can be concatenated on itself in a certain way to reduce errors arbitrarily. A directly related method also applies to over-rotational errors in two-qubit gates JonesPLA03 .Therefore this composite pulse sequence allows for an arbitrarily accurate set of universal gates, giving a threshold result for this error model.
We also consider another error model of systematic errors in even the rotation axis $`\widehat{n}`$: $`V=\mathrm{exp}(i\stackrel{}{ϵ}\stackrel{}{\sigma })`$. Here the error $`\stackrel{}{ϵ}`$ may depend on the rotation angle $`2\theta `$ but, except for a specific coordinate change, not on the axis $`\widehat{n}`$. For example, $`R_0`$ and $`R_t`$ are related by the coordinate change $`R_t=UR_0U^{}`$. We require that the errors be related by the same coordinate change, or $`\overline{R}_t=U\overline{R}_0U^{}`$.
Section II describes the basic idea behind the composite pulse sequence, by explaining its behavior when the $`\pi /3`$ rotations are perfect. In the two following sections, we extend the error model to the two cases described above.
Composite pulse sequences are an important, practical quantum control tool for removing systematic errors in a variety of quantum information processing implementations CumminsJones00 ; Riebeetal04 .We need to show that our correction sequence remains practical. While the error models we address are more general than the linear over-rotational errors which have previously been considered, the control requirement is also stricter. We typically require the ability to rotate about an arbitrary axis in the Bloch sphere, not just one in the xy plane. If rotations are only allowed about axes in the xy plane as in most NMR-type models, then our method only applies to correct $`\pi `$ pulses. In Sec. V, we compare our method, with $`\pi `$ pulses and linear over-rotational errors, to previous fully-compensating composite pulse sequences, particularly those recently discovered in BrownHarrowChuang04 .
## II Perfect $`\pi /3`$ pulses
It is instructive to start with just the core idea of our composite pulse sequence, and build up the analysis from there. Consider the sequence
$$\overline{U}^{(X)}=\overline{R}_t^{}\overline{U}\overline{R}_0\overline{U}^{}\overline{R}_t\overline{U}\overline{R}_0^{}=U\stackrel{~}{R}_0^{}V\overline{R}_0V^{}\stackrel{~}{R}_0VR_0^{},$$
(1)
where again $`VU^{}\overline{U}`$. Here we maintain a distinction between $`\stackrel{~}{R}_0U^{}\overline{R}_tU`$ and $`\overline{R}_0`$ because the errors in the two terms might be different. For the rest of this section, however, assume $`\overline{R}_0`$ and $`\overline{R}_t`$ are perfect, so $`\stackrel{~}{R}_0=\overline{R}_0=R_0\mathrm{exp}(i\frac{\pi }{6}X)`$.
Write $`V=\mathrm{exp}(i\stackrel{}{ϵ}\stackrel{}{\sigma })`$, where $`\sigma =(X,Y,Z)`$. Generally, each term of $`\stackrel{}{ϵ}`$ may be nonzero. To measure the closeness of $`\overline{U}^{(X)}`$ to $`U`$, we compute a power series expansion of $`\mathbf{(}\mathrm{Tr}(XU^{}\overline{U}^{(X)}),\mathrm{Tr}(YU^{}\overline{U}^{(X)}),\mathrm{Tr}(ZU^{}\overline{U}^{(X)})\mathbf{)}`$. We obtain
$$\begin{array}{c}\mathbf{(}2iϵ_xi\sqrt{3}(ϵ_y^2+ϵ_z^2)+O(|\stackrel{}{ϵ}|^3),\hfill \\ \hfill 2iϵ_y^3+2iϵ_yϵ_z^2+O(|\stackrel{}{ϵ}|^5),\\ \hfill 2iϵ_z^3+2iϵ_zϵ_y^2+O(|\stackrel{}{ϵ}|^5)\mathbf{)}.\end{array}$$
The first term is first order in $`ϵ_x`$ because our correction $`R_0`$ is a rotation about the $`x`$ axis, and commutes with errors in the $`X`$ direction. Errors in the $`Y`$ and $`Z`$ directions are symmetrically cancelled out, leaving only third-order terms.
We can express this result quite simply. Assume $`ϵ_{x,y,z}`$ is an $`a,b,c`$th order term. Then the $`X`$ direction error order after the $`X`$ direction composite pulse correction is applied is $`\mathrm{min}\{a,2b,2c\}`$, the $`Y`$ direction error order is $`\mathrm{min}\{3b,b+2c\}`$ and symmetrically for the $`Z`$ error. In shorthand, we write
$$(a;b;c)\underset{𝑋}{}(a,2b,2c;3b,b+2c;3c,c+2b).$$
(2)
The underset $`X`$ here refers to $`X`$ correction, and it is understood that we take a minimum on each of the three terms on the right. This notation lets us quickly understand what happens when we concatenate correction sequences. To concatenate when $`\pi /3`$ pulses are perfect, just substitute the previous level’s composite pulse sequence for $`\overline{U}`$. A level $`k`$ concatenation will require $`n_k=3n_{k1}+4`$ pulses, so the sequence length grows like $`4^k`$. For example, starting with only first-order $`Z`$ error, and applying an $`X`$ correction gives
$$(\mathrm{};\mathrm{};1)\underset{𝑋}{}(2;\mathrm{};3).$$
At this point, it is best to apply a $`Y`$ correction, since that cancels out errors in both $`X`$ and $`Z`$ directions. ($`Y`$ correction is symmetrical to $`X`$ correction, except with $`\pi /3`$ rotations about the $`y`$ axis and the same axis conjugated by $`U`$.) At the next level of concatenation, $`Z`$ correction will be optimal, and so on:
$`(\mathrm{};\mathrm{};1)`$ $`\underset{𝑋}{}`$ $`(2;\mathrm{};3)\underset{𝑌}{}(6;4;7)`$ (3)
$`\underset{𝑍}{}`$ $`(14;12;7)\underset{𝑋}{}(14;26;21)`$
$`\underset{𝑌}{}`$ $`(42;26;49)\underset{𝑍}{}(94;78;49).`$
After $`3^6=729`$ pulses of $`\overline{U}`$ or $`\overline{U}^{}`$, and $`1456`$ perfect $`\pi /3`$ pulses (about six axes), the error is only $`O(|\stackrel{}{ϵ}|^{49})`$.
###### Remark 1 (Generalization).
The question of whether this pulse sequence generalizes deserves further study. We have been able to find a pulse sequence with five applications of $`\overline{U}`$ or $`\overline{U}^{}`$, and perfect rotations by $`\pi /5`$ or $`3\pi /5`$, which on input $`|\pm \frac{1}{\sqrt{2}}(|0\pm |1)`$ achieves a fidelity error of $`O(ϵ^{10})`$:
$$\begin{array}{c}|\pm |U^{}\overline{U}R_0(\frac{3\pi }{5})\overline{U}^{}R_t(\frac{\pi }{5})\overline{U}R_0(\frac{\pi }{5})\overline{U}^{}R_t(\frac{3\pi }{5})\overline{U}|\pm |^2\hfill \\ \hfill =1O(ϵ^{10}).\end{array}$$
(4)
(See Fig. 1 for geometrical intuition.) However, this sequence gives no improvement with imperfect correction rotations.
###### Remark 2 (Error measurement).
For us it is key to measure the direction of the error, as well as its magnitude. How does our method of measuring error compare to other reasonable methods? On a particular input state, the difference in the fidelity from one is quadratically smaller than our measure. The so-called infidelity between $`U`$ and $`\overline{U}^{}`$, or $`1\frac{1}{2}|\mathrm{Tr}U^{}\overline{U}^{}|`$ is used in McHughTwamley04 ; CumminsJones00 ; JonesPLA03 ,and is also quadratically smaller than our measure. Brown et al BrownHarrowChuang04 use as their measure of distance the trace distance $`\mathrm{Tr}|U\overline{U}^{}|`$, which depends on the global phase of the operators. Our correction sequence does not give higher order accuracy in the global phase, but a simple modification of the trace distance optimizes over global phases, and then this measure of error is of the same order as ours.
## III Imperfect $`\pi /3`$ pulses, error angle-dependent & axis-independent
Let us consider the more realistic case that the $`\pi /3`$ pulses are themselves erroneous. Assume that the error in a rotation depends only on the rotation angle, not on the rotation axis. When we attempt to apply $`\mathrm{exp}(i\theta \widehat{n}\stackrel{}{\sigma })`$, we actually apply $`\mathrm{exp}\{i[\theta +ϵ(\theta )]\widehat{n}\stackrel{}{\sigma }\}`$. Here $`ϵ(\theta )`$ can be an arbitrary function of $`\theta `$, which is however always small \[order $`O(|ϵ|)`$\]. The error amount does not depend on the rotation axis $`\widehat{n}`$. Previous work has only considered the less-general case of linear errors, $`ϵ(\theta )=ϵ\theta `$.
In fact, let us generalize our calculations slightly further. We will allow errors in the $`\pi /3`$ pulses besides just over-rotation errors, except these errors must be the same under change of coordinates by $`U`$. In particular, $`\overline{R}_t=U\overline{R}_0U^{}`$. It isn’t clear in what physical models this will be an appropriate base error assumption – perhaps one in which the entire apparatus for applying a $`\pi /3`$ rotation is rotated about the qubit in three dimensions as $`U`$ acts in $`𝐑^\mathrm{𝟑}`$, or equivalently, the qubit is physically rotated. However, the added generality will be necessary for considering concatenation of this correction sequence.
Write $`\overline{R}_0=R_0\mathrm{exp}(i\stackrel{}{\delta }\stackrel{}{\sigma })`$. We obtain
$`\mathrm{Tr}(XU^{}\overline{U}^{(X)})`$ $`=`$ $`2iϵ_x+2i(\sqrt{3}\delta _y+\delta _z)ϵ_yi\sqrt{3}(ϵ_y^2+ϵ_z^2)2i(\delta _y\sqrt{3}\delta _z)ϵ_z+O(|\stackrel{}{ϵ}|^3+|\stackrel{}{\delta }||\stackrel{}{ϵ}|^2+|\stackrel{}{\delta }|^2|\stackrel{}{ϵ}|)`$
$`\mathrm{Tr}(YU^{}\overline{U}^{(X)})`$ $`=`$ $`\left(\begin{array}{c}\hfill 2i(\sqrt{3}\delta _y+\delta _z)ϵ_x4i(\sqrt{3}\delta _x\sqrt{3}\delta _y\delta _z+\delta _z^2)ϵ_y\sqrt{3}i(\sqrt{3}\delta _y+\delta _z)ϵ_y^2+2iϵ_y^3\\ \hfill 2i(\sqrt{3}\delta _y^22\delta _y\delta _z\sqrt{3}\delta _z^2)ϵ_z\sqrt{3}i(\sqrt{3}\delta _y\delta _z)ϵ_z^2+2iϵ_yϵ_z^2+O(|\stackrel{}{ϵ}|^3+|\stackrel{}{\delta }|^2|\stackrel{}{ϵ}|)\end{array}\right).`$
$`\mathrm{Tr}(ZU^{}\overline{U}^{(X)})`$ can be determined by symmetry. In our shorthand notation, with $`\delta _{x,y,z}`$ being $`d,e,f`$th order, respectively,
$$(a;b;c)\underset{𝑋}{}\left(\begin{array}{c}a,e+b,f+b,2b,e+c,f+c,2c;\\ e+a,f+a,d+b,e+f+b,2f+b,e+2b,f+2b,3b,2e+c,e+f+c,2f+c,e+2c,f+2c,b+2c;\\ e+a,f+a,2e+b,e+f+b,2f+b,e+2b,f+2b,d+c,2e+c,e+f+c,2b+c,e+2c,f+2c,3c\end{array}\right).$$
(5)
For example, taking $`d=e=f=\mathrm{}`$, we recover Eq. (2) from the perfect $`\pi /3`$ pulse case. In the case of first-order over-rotation, $`d=1`$, $`e=f=\mathrm{}`$,
$$(\mathrm{};\mathrm{};1)\underset{X,Y,Y}{}(4;4;4).$$
(6)
To obtain arbitrarily accurate rotations, it is most effective to correct both the applications of $`\overline{U}`$ and the $`\pi /3`$ correction pulses. So at this point correct the $`\pi /3`$ pulses until $`d=e=f=4`$. Note that applying such a correction maintains the invariant that the error in a pulse depend only on the angle and not the axis. Now
$$(4;4;4)\underset{X,Y,Z}{}(12;12;12).$$
(7)
Every three levels of concatenation (both on $`\overline{U}`$ and the $`\pi /3`$ pulses) increases the error order by a factor of three. Therefore obtaining error tolerance to a desired amount $`ϵ_{}`$ requires poly-logarithmically many pulses in $`1/ϵ_{}`$.
## IV Imperfect $`\pi /3`$ pulses, error both angle- and axis-dependent
What if the error in $`\pi /3`$ pulses depends on which basis they are carried out in, i.e., $`\overline{R}_tU\overline{R}_0U^{}`$? Can we still obtain arbitrarily accurate pulses? Perhaps surprisingly, the answer is yes, if the error is of a restricted form: only over-rotation errors. However, the orders will not grow exponentially quickly in the number of concatenation levels, only linearly, implying that error tolerance to an amount $`ϵ_{}`$ will require polynomially many pulses in $`1/ϵ_{}`$ instead of only poly-logarithmically many.
Write $`\overline{R}_0=R_0\mathrm{exp}(i\delta X)`$, $`\overline{R}_t=UR_0\mathrm{exp}(i\widehat{\delta }X)U^{}`$. Expanding $`\mathbf{(}\mathrm{Tr}(XU^{}\overline{U}^{(X)}),\mathrm{Tr}(YU^{}\overline{U}^{(X)}),\mathrm{Tr}(ZU^{}\overline{U}^{(X)})\mathbf{)}`$, we obtain
$$\left(\begin{array}{c}2iϵ_xi\sqrt{3}\left(ϵ_y^2+ϵ_z^2\right)+O\left(\left|\delta \right|ϵ^2+\left|\stackrel{}{ϵ}\right|^3\right),\\ 2i\sqrt{3}\left(\delta +\widehat{\delta }\right)ϵ_y2i\left(\delta \widehat{\delta }\right)ϵ_z+2iϵ_y^3+2iϵ_yϵ_z^2+O\left(\delta ^2\left|\stackrel{}{ϵ}\right|+\left|\stackrel{}{ϵ}\right|^5\right),\\ 2i\sqrt{3}\left(\delta +\widehat{\delta }\right)ϵ_z2i\left(\delta \widehat{\delta }\right)ϵ_y+2iϵ_z^3+2iϵ_zϵ_y^2+O\left(\delta ^2\left|\stackrel{}{ϵ}\right|+\left|\stackrel{}{ϵ}\right|^5\right)\end{array}\right),$$
assuming $`|\widehat{\delta }|=\mathrm{\Theta }(|\delta |)`$. In our shorthand notation, with $`\delta `$ first order, the rule is
$$\begin{array}{c}(a;b;c)\underset{𝑋}{}\hfill \\ \hfill (a,2b,2c;3b,b+2c,1+b,1+c;3c,c+2b,1+b,1+c).\end{array}$$
(8)
For example,
$`(\mathrm{};\mathrm{};1)`$ $`\underset{𝑋}{}`$ $`(2;2;2)\underset{𝑋}{}(2;3;3)`$ (9)
$`\underset{𝑌}{}`$ $`(3;3;3)\underset{𝑋}{}(3;4;4),`$
and so on, with every two levels of concatenation increasing the error order by one. Note that we do not concatenate corrections onto the $`\pi /3`$ pulses, because then the error would no longer be simply over-rotational. (Even with a more general expansion, it turns out that the convergence is still only be linear in the number of concatenation levels.)
## V $`\pi `$ pulses in NMR
While our method corrects against more general types of errors than previous composite pulse sequences, it also has a stronger requirement. Namely, we must be able to apply a $`\pi /3`$ rotation about the $`x`$, $`y`$ and $`z`$ axes, and also about those same axes in the $`U`$-transformed basis. In most current proposed quantum information implementations, primitive rotations are only allowed about axes in the $`xy`$ plane. For $`U`$, $`R_t`$ and $`R_0`$ all to be rotations about axes in the $`xy`$ plane, it must be that $`U`$ is a rotation by an integer multiple of $`\pi `$. This is a considerable restriction on the applicability of our method. Still, for $`\pi `$ pulses our correction succeeds in a setting more general than that for which previous methods could correct; for example, we can correct for non-linear over-rotations.
Assume now that the systematic error is in fact a linear over-rotation; when we try to apply $`\mathrm{exp}(i\theta \widehat{n}\stackrel{}{\sigma })`$, we actually apply $`\mathrm{exp}[i\theta (1+ϵ)\widehat{n}\stackrel{}{\sigma }]`$. Assume $`U=\mathrm{exp}(i\frac{\pi }{2}X)=X`$. Our $`\pi /3`$ correction method, correcting in the $`Y`$ direction, leaves behind second-order errors, with a sequence length of $`3\pi +4\pi /3`$. (Practical composite pulse sequences need to be as short as possible, in order to minimize any non-ideal effects not accounted for in our error model.) We cannot concatenate a $`Z`$ correction, but can concatenate alternately $`X`$ and $`Y`$ corrections to reduce the error arbitrarily. Concatenating an $`X`$ correction onto the $`Y`$ correction leaves behind a third-order error, with a sequence length of $`(14\frac{1}{3})\pi `$.
How does the $`\pi /3`$ correction method compare with previous correction methods? The most practical correction methods previously known were the B2 (also known as BB1) correction sequence, and the recently discovered B4 sequence BrownHarrowChuang04 . These sequences are implemented as follows:
B2: $`\left(\overline{R}_{\varphi _1}(\pi )\overline{R}_{3\varphi _1}(2\pi )\overline{R}_{\varphi _1}(\pi )\right)\overline{U}`$ (10)
B4: $`\left(\begin{array}{c}\hfill (\overline{R}_{\varphi _2}(\pi )\overline{R}_{3\varphi _2}(2\pi )\overline{R}_{\varphi _2}(\pi ))^4\\ \hfill \times (\overline{R}_{\varphi _2}(2\pi )\overline{R}_{\varphi _2}(4\pi )\overline{R}_{\varphi _2}(2\pi ))\\ \hfill \times (\overline{R}_{\varphi _2}(\pi )\overline{R}_{3\varphi _2}(2\pi )\overline{R}_{\varphi _2}(\pi ))^4\end{array}\right)\overline{U},`$ (11)
where $`\overline{R}_\varphi (\theta )\mathrm{exp}(i\frac{\theta (1+ϵ)}{2}(\mathrm{cos}\varphi X+\mathrm{sin}\varphi Y))`$, $`\mathrm{cos}\varphi _1=\frac{\theta }{4\pi }`$, and $`\mathrm{cos}\varphi _2=\frac{\theta }{24\pi }`$. They leave behind third- and fifth-order errors, respectively. The total sequence length for correcting a $`\pi `$ rotation is $`5\pi `$ for B2 and $`41\pi `$ for B4. Hence our method, in this linear over-rotation error model, seems to offer little over the plain B2 sequence. Table 1 gives values for the infidelities $`1\frac{1}{2}|\mathrm{Tr}U^{}\overline{U}^{()}|`$ for various correction sequences. (As previously remarked, the infidelity is quadratically smaller than the trace distance, so the infidelity error orders for the three possibilities $`\pi /3`$ $`Y`$, B2 and B4 are $`4`$, $`6`$ and $`10`$, respectively.)
Does our method complement previous correction methods? To answer this question, we must determine the direction of the error left behind after a correction sequence. For example, B2 and B4 are each implemented as $`Bi\overline{U}`$, where $`Bi`$ is some particular pulse sequence not involving $`\overline{U}`$. A simple calculation shows that both B2 and B4 leave behind an error which is has relatively large $`X`$ and $`Y`$ components. Therefore, concatenating on $`X`$ or $`Y`$ correction will not increase the error order. We can however find an axis in the $`xy`$ plane which is approximately orthogonal to the $`xy`$ component of the error, and correct along this axis. Alternatively, we can symmetrize the B2 and B4 sequences into $`\mathrm{exp}[i\frac{\pi }{4}(1+ϵ)X]Bi\mathrm{exp}[i\frac{\pi }{4}(1+ϵ)X]`$. In this more symmetrical form, the error magnitude is unchanged, but the direction is entirely into the $`xz`$ plane. Therefore, we can simply apply a $`Y`$ correction to the symmetrized sequences. Table 1 compares the infidelities of $`Y`$ correction concatenated onto the symmetrized B2 and B4 correction sequences. Note that the former case gives fourth-order protection with a sequence length of only $`(16\frac{1}{3})\pi `$; this gives a new, perhaps practical, compromise between the B2 and B4 correction sequences. Figure 2 plots the fidelities for $`ϵ>0`$.
## VI Conclusion
We have presented two main results. The $`\pi /3`$ correction sequence protects against general errors which depend arbitrarily on the rotation angle but not the rotation axis.The same sequence protects against over-rotational error which depends arbitrarily on both the rotation axis and the angle of rotation.Previously, composite pulse correction sequences were only known for the cases when the error was independent of the rotation axis, and depended linearly on the rotation angle.
Moreover, our composite pulse correction sequence concatenates nicely to reduce errors arbitrarily. In the first case, the overhead number of pulses is poly-logarithmic in the desired accuracy, and in the second case the overhead is polynomial.
However, our correction sequence in general requires primitive rotations about arbitrary axes in the Bloch sphere, and only applies to correct $`\pi `$ pulses in the typical situation in which rotations are only allowed about axes in the xy plane. For correcting $`\pi `$ pulses, our method concatenated on top of a B2 pulse correction provides a new compromise between B2 and B4.
B. R. acknowledges support from NSF ITR Grant CCR-0121555, and ARO Grant DAAD 19-03-1-0082.
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# Superfluid-Insulator Transitions on the Triangular Lattice
## I Introduction
Theoretical interest in quantum phase transitions from superfluid to Mott insulating states of bosons has recently been revived by their experimental realization in cold atoms trapped in an optical lattice.Greiner Extensions of these experiments are expected to soon provide a great variety of real life toy models, where theoretical scenarios for such phase transitions can be tested.Demler03 ; Zoller04 ; Zoller05 For the condensed matter community, such transitions are interesting in their own right, but also provide a simpler context in which some aspects of Mott conducting-insulating transitions of electrons can be explored. A better understanding of such Mott criticality generally may help to explain mysteries in various strongly correlated materials, from heavy fermion metals livrefs to cuprate superconductors,cuprefs in which Mott criticality may plausibly be argued to play a key role.
An exciting theoretical development in the field has been the discovery that some quantum phase transitions require a fundamentally new description, not based on the now-standard Landau’s concept of an order parameter. Instead, such quantum critical points (QCPs) are described in terms of emergent degrees of freedom, not present in either phase and appearing due to certain special dynamically generated low-energy symmetries at the critical point.dcprefs An interesting consequence of the emergent low-energy symmetry of the critical point is the near-degeneracy of “competing ordered” states unrelated by any microscopic symmetry (but unified with one another by the emergent one) in the neighborhood of the quantum phase transition. There is, at present, unfortunately, no general way to a priori identify the appropriate emergent degrees of freedom, should they exist, for any putative quantum critical point.
In the particular context of two-dimensional bosonic superfluid-insulator transitions, a general non-Landau-Ginzburg-Wilson (non-LGW) framework has recently been proposed in Ref. Balents04, (and see Ref. ykis, for a pedagogical review), and carried out explicitly on the square lattice. In particular, the Mott transition can be described in terms of the vortex excitations of the superfluid. In a two dimensional superfluid, vortices are point-like “particles” whose creation/annihilation operators can be used to construct a quantum field theory. The vortices being non-local topological objects, these vortex fields are, however, not themselves order parameters in the LGW sense. The non-locality is manifested by the presence of a (non-compact) $`U(1)`$ gauge field, to which the vortex fields are coupled in the “dual” vortex field theory. This formulation is general because it is based on the excitations of the superfluid, which is a stable and apparently featureless (i.e. without broken symmetry apart from off-diagonal long range order) state at any boson density. Nevertheless, it was shown in Ref. Balents04, that the vortices exhibit a subtle quantum order which is sensitive to the boson density $`f`$ (per unit cell of the lattice). In particular, at non-integral $`f`$, the vortices form non-trivial multiplets transforming under a projective symmetry group (PSG - technically, a projective representation of the lattice space group). The Lagrange density for the vortex field theory therefore takes the general form
$``$ $`=`$ $`{\displaystyle \underset{\mathrm{}=0}{\overset{N}{}}}\left[|(_\mu iA_\mu )\phi _{\mathrm{}}|^2+s|\phi _{\mathrm{}}|^2\right]+_{\mathrm{int}}[\{\phi _{\mathrm{}}\}]`$ (1)
$`+{\displaystyle \frac{1}{2e^2}}(ϵ_{\mu \nu \lambda }_\nu A_\lambda )^2,`$
where $`\phi _{\mathrm{}}`$ with $`\mathrm{}=1\mathrm{}N`$ are vortex fields for the $`N`$ members of the multiplet ($`N`$ depends upon $`f`$ – see Sec. II), and $`A_\mu `$ is the dual $`U(1)`$ gauge field. Quartic and higher order terms are contained in $`_{\mathrm{int}}`$, the structure of which is dictated by the PSG. The Mott transition is captured by Eq.(1) in a simple way. The ground state of $`s>0`$ corresponds to the vacuum of vortices – flux is expelled from the system, so this is a superfluid. In the $`s<0`$ phase, at least one of the vortex flavors will condense, and the gauge fields will acquire a gap by the Higgs mechanism. This is indicative of a charge gap, and describes an incompressible Mott insulating phase.
Transformations within the $`\phi _{\mathrm{}}`$ multiplet comprise the emergent symmetry operations of the “deconfined” quantum critical point. This structure has an important physical consequence as well: it naturally and unavoidably leads to (particular) broken spatial symmetries in the Mott state. A mean-field analysis of the vortex field theory predicts a direct superfluid to Mott transition, as well as the nature of the charge ordering in the Mott phase.
In this paper, we extend this vortex field theory approach to describe superfluid to Mott insulator transitions of bosons on the triangular lattice (hexagonal lattice in proper crystallographic nomenclature) at fractional boson fillings $`f=p/q`$, with $`p,q`$ relatively prime integers. We focus particularly on the most interesting examples, $`f=1/2`$ and $`f=1/3`$. The $`1/2`$-filling case introduces a new ingredient not present on the square lattice: geometrical frustration. Here we refer to “charge frustration” of the ordering of localized boson configurations in the presence of short-range repulsive interactions. A consequence – or perhaps definition – of such geometrical frustration is the exact or near degeneracy of many distinct ordered states. The similarity of this property with the emergent near-degeneracy of competing orders near a deconfined quantum critical point suggests a possible link between the two phenomena. This connection indeed appears to be borne out by the analysis in this paper.
In the simplest classical models of frustration, the degeneracy amongst low-energy states is not only large but macroscopic (i.e. with an entropy proportional to the sample volume). This classical degeneracy, when lifted by quantum fluctuations, may produce unusual ground states. A number of very recent papersMelko05 ; Damle05 ; Troyer05 have investigated the system of hardcore bosons with nearest-neighbor interactions (which can be mapped to a spin-$`1/2`$ XXZ model) on the triangular lattice near half-filling. This realizes such approximately classical frustration in the limit of very strong near-neighbor repulsion. These works demonstrated that in this system the lifting of the macroscopic classical degeneracy results in an unusual supersolid ground state, which we denote SS3 because of its 3-sublattice structure. This “order by disorder” mechanism is very different from the “conventional” (theoretically!) picture of supersolidity, via a condensation of vacancies and/or interstitials in an ordered solid.supersolid
Since the above studies clarified that such Mott states do not occur in the simplest nearest-neighbor interaction model, it is apparent that microscopically, longer-range interactions, possibly including ring-exchange,Sandvik02 are necessary to observe these transitions in a microscopic model. Finding simple interactions that produce nontrivial insulating ground states on the triangular lattice is an important and difficult problem, that will likely require sophisticated numerical methods. We will not address this issue here (but see the Discussion, Sec. V).
Our phenomenological vortex field theory, on the contrary, describes universal aspects (independent of microscopic realization) of Mott insulating, superfluid, and other states, and the transitions between them. As for the square lattice case, a mean-field analysis predicts a direct superfluid-Mott transition, with a diverse set of Mott insulating phases. Also like the square lattice, the vortex field theory has an enhanced emergent symmetry at the critical point, a hallmark of deconfined criticality. A significant difference from these prior examples of deconfined criticality, however, is that, in the frustrated case, $`f=1/2`$, the emergent symmetry is nonabelian, containing an $`SU(2)`$ “pseudospin” subgroup. The much larger (than in non-frustrated cases) emergent symmetry can be understood physically as symptomatic of the larger near-degeneracies present in this case due to frustration.
Remarkably, going beyond the mean-field analysis of the vortex field theory, our approach connects very nicely to the supersolid phase of the XXZ model. Indeed, the supersolid order parameter – describing the growth of the supersolid state out of the featureless superfluid – appears in a particularly simple form in the vortex variables. Moreover, our approach reveals an alternative view of the supersolid, as a partially-melted “parent” Mott insulating state, with “quantum disordered” pseudospin. This loss of pseudospin order simultaneously with the onset of superfluidity is possible because, as we show, the pseudospin skyrmion excitation of the Mott insulator carries physical boson charge. The supersolid may thereby also be viewed as a condensate of these skyrmions. Furthermore, this picture leads directly to the prediction that the Mott insulator to SS3 transition is described by the recently discovered Non-Compact CP<sup>1</sup> (NCCP<sup>1</sup>) quantum critical universality class.Lesik This transition describes the quantum disordering of a pseudospin vector $`𝐒`$ (the order parameter for the additional solid order of the Mott state) in $`2+1`$ dimensions, when “hedgehog” instantons are absent in space-time. These hedgehog events correspond precisely to processes which change the skyrmion number, and are therefore prohibited by charge conservation in this case.
The paper is organized as follows. In Sec.II we develop the dual vortex theory for the triangular lattice at a general rational filling. We derive the PSG transformations of the degenerate low energy vortex modes and discuss some of their general properties. In Sec. III we apply this theory to the two cases $`f=1/3`$ and $`f=1/2`$, and discuss the resulting Mott states that are obtained by a mean field analysis of the vortex field theory. In Sec. IV we present a “hard-spin” formulation of the vortex action, which enables a study of its phases and transitions beyond mean field theory. We discuss the new phases which arise, notably supersolids, the elementary excitations of the different states, and the transitions amongst them. This includes notably the SS3 phases and the NCCP<sup>1</sup> transitions between them and their parent Mott insulator. We conclude with a discussion in Sec. V of a more microscopic physical picture for the most interesting Mott and supersolid states at half-filling, the connection to the recent studies of the XXZ model, and the prospects for observing more of the physics in this paper in related models. A number of appendices provide useful details of various technical results.
## II Continuum dual vortex theory
As discussed in the introduction, our aim is to derive a field theoretic description of the vicinity of a superfluid to Mott insulator transition, in terms of the vortex excitations of the superfluid state. As is well-known, lattice models of bosons can indeed be reformulated in vortex variables on the dual lattice, a technique called duality. As detailed for instance in Ref. ykis, (and references therein), this can in principle be carried out exactly for any lattice boson Hamiltonian, e.g. Bose-Hubbard models, XXZ models, etc. Unfortunately, this exact mapping does not lead to a particularly tractable limit of the vortex theory, and it is therefore difficult to extract quantitative predictions directly from the lattice vortex theory.
Fortunately, we can avoid this difficulty in addressing our stated goal of understanding universal phenomena in the vicinity of the superfluid-Mott quantum critical point. For this purpose, we need not specify any particular microscopic boson model. Instead, we will illustrate the calculations through the use of a very simple lattice vortex theory, which is chosen to have the same spatial symmetries as the physical triangular lattice, and to exhibit a superfluid to Mott insulation transition. The universal properties of interest will coincide with those of more physical microscopic boson models.
In the Euclidean vortex coherent-state path integral formulation in which we work, the lattice vortex action is
$`S`$ $`=`$ $`t_v{\displaystyle \underset{a\mu }{}}[\psi _a^{}e^{iA_{a\mu }}\psi _{a+\mu }+c.c.]`$ (2)
$`+`$ $`{\displaystyle \underset{a}{}}[s|\psi _a|^2+u|\psi _a|^4]`$
$`+`$ $`{\displaystyle \frac{1}{2e^2}}{\displaystyle \left(ϵ_{\mu \nu \lambda }\mathrm{\Delta }_\nu A_{a\lambda }2\pi f\delta _{\mu \tau }\right)^2}.`$
Here $`a`$ labels sites of the dual 2+1-dimensional uniformly stacked honeycomb lattice, with the vertical direction being the imaginary time direction $`\tau `$, and $`\mu `$ is summed over nearest-neighbor links (3 at $`120`$ degrees connecting spatial neighbors and a fourth along the imaginary time direction). The action is written in terms of the complex vortex field $`\psi _a`$ (and its conjugate $`\psi _a^{}`$) which annihilates (creates) unit vorticity, as well as a dual gauge field $`A_{a\mu }`$. The physical meaning of the gauge field $`A_{a\mu }`$ is that its curl gives ($`2\pi `$ times) the bosonic 3-current density, and importantly, the temporal component of the current density is ($`2\pi `$ times) the charge density. For a bosonic system with density $`f`$ bosons per site, we must therefore enforce the condition that $`2\pi f`$ flux on average passes through each hexagonal plaquette of the honeycomb lattice. We will assume the filling to be rational $`f=p/q`$, where $`p`$ and $`q`$ are relatively prime integers, and will mostly concentrate on the two cases $`f=1/2,1/3`$. Physically, $`t_v`$ represents a vortex hopping amplitude, $`s`$ and $`u`$ represent short-range vortex “core” energies and interactions, and $`e^2`$ represents the strength of the dual electromagnetic field fluctuations (it is roughly proportional to the local superfluid density away from the Mott quantum critical point).
We may think of the parameter $`s`$ in Eq.(2) as driving the superfluid-insulator transition. Large positive $`s`$ corresponds to the superfluid phase, in which vortices are gapped, and large negative $`s`$ corresponds to the set of possible insulating phases, in which vortices are condensed. It is convenient to think about the transition in Eq. (2) for small $`e^2`$, i.e. neglecting to a first approximation dual gauge fluctuations – they will however be restored at a later stage of analysis. In this limit we can treat the dual gauge field in a mean field approximation and take:
$$ϵ_{\mu \nu \lambda }\mathrm{\Delta }_\nu A_{a\lambda }=2\pi f\delta _{\mu \tau }.$$
(3)
As $`s`$ is decreased from large positive values, we expect an instability of the “vortex vacuum” when the energy of the lowest vortex excitation approaches zero – or equivalently, when the minimum eigenvalue of the quadratic form for $`\psi _a,\psi _a^{}`$ in Eq. (2) vanishes. To study the universal critical properties of the superfluid-insulator transition, it is sufficient to isolate these low-energy particles, which comprise the vortex multiplets discussed in the introduction. The continuum limit of Eq. (2) then consists of a set of vortex fields, each representing one of these minimum energy vortex particles. We will also take the trivial continuum limit of Eq.(2) in the temporal direction.
Because of the non-zero gauge flux through each spatial plaquette, the minimum energy multiplets are non-trivial. The form of the continuum action in this case is determined by the projective representation of the space group (PSG), Balents04 under which the vortex fields $`\psi `$ transform. To find it, we must work through the consequences of some specific gauge choice. As in Ref.Balents04, , we will choose the Landau gauge for $`A_{a\mu }`$. Namely, let
$$𝐚_1=\widehat{x},𝐚_2=\frac{1}{2}\widehat{x}+\frac{\sqrt{3}}{2}\widehat{y},$$
(4)
be the two basis vectors of the honeycomb lattice, as shown in Fig.1 (note that the honeycomb lattice has two sites per unit cell). Coordinates will be specified, when explicit, in this basis, $`𝐫=a_1𝐚_1+a_2𝐚_2`$, with integer $`a_1,a_2`$, and $`a_1=a_2=0`$ corresponding to a “type 1” site (see Fig. 1) of the dual honeycomb lattice. Landau gauge is
$$A_{ay}=2\pi fa_1,$$
(5)
and $`A_{a\mu }=0`$ for all other directions; that is, only the gauge field on vertical links is chosen non-zero. The full PSG is generated by a set of unitary transformations of $`\psi `$, one for each generator of the lattice space group. For the triangular lattice, we choose the space group generators as two elementary lattice translations $`T_1,T_2`$, a $`2\pi /3`$ rotation with respect to site $`1`$ in Fig. 1, $`R_{2\pi /3}`$, and two reflections, $`I_{d_1},I_{d_2}`$.<sup>1</sup><sup>1</sup>1It is straightforward to show that the $`6`$-fold rotation about a direct lattice site, $`R_6`$, is determined from these generators by the relation $`R_6T_2R_{2\pi /3}I_{d_1}I_{d_2}=1`$, and so is not independent. The PSG transformations of the vortex fields, characterized by the unimodular complex number $`\omega =e^{2\pi if}`$, are then:
$`T_1`$ $`:`$ $`\{\begin{array}{ccc}\psi _1(a_1,a_2)\hfill & & \psi _1(a_11,a_2)\omega ^{a_2}\hfill \\ \psi _2(a_1,a_2)\hfill & & \psi _2(a_11,a_2)\omega ^{a_1+1}\hfill \end{array},`$ (8)
$`T_2`$ $`:`$ $`\{\begin{array}{ccc}\psi _1(a_1,a_2)\hfill & & \psi _1(a_1,a_21)\hfill \\ \psi _2(a_1,a_2)\hfill & & \psi _2(a_1,a_21)\hfill \end{array},`$ (11)
$`R_{2\pi /3}`$ $`:`$ $`\{\begin{array}{ccc}\psi _1(a_1,a_2)\hfill & & \psi _1(a_2a_1,a_1)\omega ^{\frac{a_1(2a_2a_11)}{2}}\hfill \\ \psi _2(a_1,a_2)\hfill & & \psi _2(a_2a_1,a_11)\omega ^{\frac{a_1(2a_2a_1+1)}{2}}\hfill \end{array},`$ (14)
$`I_{d_1}`$ $`:`$ $`\{\begin{array}{ccc}\psi _1(a_1,a_2)\hfill & & \psi _1^{}(a_2,a_1)\omega ^{a_1a_2}\hfill \\ \psi _2(a_1,a_2)\hfill & & \psi _2^{}(a_21,a_11)\omega ^{a_1(a_2+1)}\hfill \end{array},`$ (17)
$`I_{d_2}`$ $`:`$ $`\{\begin{array}{ccc}\psi _1(a_1,a_2)\hfill & & \psi _2^{}(a_2,a_11)\omega ^{a_1a_2}\hfill \\ \psi _2(a_1,a_2)\hfill & & \psi _1^{}(a_2+1,a_1)\omega ^{a_1(a_2+1)}\hfill \end{array}.`$ (20)
For the special case $`f=1/2`$, in which we are principally interested, one usually considers in addition to these spatial symmetries, an extra particle-hole symmetry, which we denote $`C`$ for “charge conjugation”. This interchanges singly-occupied and empty sites on the direct lattice. In the spin language appropriate for hard-core bosons, this is just a $`180^{}`$ rotation around the $`x`$ axis in spin space, which has the effect of an Ising transformation $`S_i^zS_i^z`$ (and likewise for $`S_i^y`$). The XXZ model, and indeed any hard core boson model at $`f=1/2`$ with only pairwise interactions, possesses such a particle-hole symmetry. In the dual theory, this requires the action to be invariant under
$$C:\psi _i(a_1,a_2)\psi _i^{}(a_1,a_2),$$
(21)
and simultaneous sign change of the fluctuating part of the dual gauge field $`A_\mu A_\mu `$.
The quadratic action of Eq. (2) in Landau gauge has a periodicity in real space of $`q`$ unit cells in the $`𝐚_1`$ direction, and one unit in the $`𝐚_2`$ direction (note that, of course, the physics itself has the full periodicity of the honeycomb (or underlying triangular direct) lattice). The eigenstates of the quadratic action can therefore be characterized by their quasimomenta in the corresponding reduced Brillouin zone. Specifically, we introduce the basis vectors of the reciprocal lattice,
$$𝐛_1=\widehat{x}+\frac{1}{\sqrt{3}}\widehat{y},𝐛_2=\frac{2}{\sqrt{3}}\widehat{y}.$$
(22)
Wavevectors will, when necessary, be specified by coordinates $`(k_1,k_2)`$, with $`𝐤=k_1𝐛_1+k_2𝐛_\mathrm{𝟐}`$ (reciprocal lattice vectors correspond to $`k_1,k_2`$ being integral multiples of $`2\pi `$).
By applying the methods of Ref.Balents04, , one may readily find the minimum energy multiplet for arbitrary $`p,q`$, and their PSG transformations. Briefly, this is accomplished by Fourier transforming Eqs. (8) to obtain the PSG for general $`\psi (𝐤)`$ in momentum space – this is given in Appendix A – and using the non-commutative algebra of translations and rotations implicit in Eqs. (8) to generate a full set of eigenfunctions. One finds two cases. For $`q`$ odd, there are $`q`$ minima of the vortex dispersion, i.e. $`N=q`$ in Eq. (1). These occur at momenta
$$𝐤_{\mathrm{}}=(0,2\pi f\mathrm{}),\mathrm{}=0,\mathrm{},q1,$$
(23)
taking wavevectors $`𝐤_{\mathrm{}}`$ to lie in the reduced Brilloin zone $`\pi /qk_1<\pi /q`$ and $`\pi k_2<\pi `$. By contrast, for $`q`$ even, $`N=2q`$, and it is convenient to divide minima into two sets $`\alpha =\pm 1\pm `$, parameterizing $`\mathrm{}=(1+\alpha )q/2+\sigma `$, $`\sigma =0\mathrm{}q1`$. The vortex field operator $`\phi _{\alpha \sigma }`$ then acts on eigenstates with the wavevector
$$𝐤_{\alpha \sigma }=(\pi \alpha /3q,\pi \alpha /3q+2\pi f\sigma )$$
(24)
in the reduced Brillouin zone. For even $`q`$, therefore, Eq. (1) may be rewritten as
$`_0^{even}`$ $`=`$ $`{\displaystyle \underset{\alpha =\pm }{}}{\displaystyle \underset{\sigma =0}{\overset{q1}{}}}\left[|(_\mu iA_\mu )\phi _{\alpha \sigma }|^2+s|\phi _{\alpha \sigma }|^2\right]`$ (25)
$`+`$ $`{\displaystyle \frac{1}{2e^2}}(ϵ_{\mu \nu \lambda }_\nu A_\lambda )^2.`$
We note in passing that the set of $`\phi _{\mathrm{}}`$ corresponding to these wavevectors gives the smallest number of minima possible for the vortex dispersion (i.e. they comprise the smallest-dimensional representation of the PSG for the given $`f`$), and they are what is realized without special fine-tuning of the vortex kinetic energy terms.
The PSG transformations of the multiplet can be found using Eq.(91) of Appendix A, by a straightforward generalization of the procedure, detailed in Ref. Balents04, . For instance, under the translations, one finds
$`T_1:`$ $`\phi _{\mathrm{}}\phi _\mathrm{}1,`$
$`T_2:`$ $`\phi _{\mathrm{}}\omega ^{\mathrm{}}\phi _{\mathrm{}},`$ (26)
for $`q`$ odd, and
$`T_1:`$ $`\phi _{\alpha \sigma }e^{i\pi \alpha /3q}\phi _{\alpha ,\sigma 1},`$
$`T_2:`$ $`\phi _{\alpha \sigma }e^{i\pi \alpha /3q}\omega ^\sigma \phi _{\alpha \sigma },`$ (27)
for $`q`$ even. Here and in the following, the index $`\mathrm{}`$ for odd $`q`$ and $`\sigma `$ for even $`q`$ will be regarded as cyclic modulo $`q`$. The remaing PSG generators are given in Eqs. (A,A) in Appendix A.
It is now straightforward to write down the most general continuum vortex Lagrangian density, describing the superfluid-insulator transition on the triangular lattice. The most general terms at quadratic order are simply given by Eqs. (1,25). Let us now consider the quartic potential terms in the continuum theory. Following the general approach of Ref. Balents04, , we first write down the continuum Lagrangian, imposing only the restrictions from gauge symmetry and translational symmetry, i.e. invariance under Eqs. (II,II). One finds
$`_1^{\mathrm{odd}}`$ $`=`$ $`{\displaystyle \underset{\mathrm{}mn}{\overset{q1}{}}}\gamma _{mn}\phi _{\mathrm{}}^{}\phi _{\mathrm{}+m}^{}\phi _{\mathrm{}+n}^{}\phi _{\mathrm{}+mn}^{},`$ (28)
$`_1^{\mathrm{even}}`$ $`=`$ $`{\displaystyle \underset{\sigma ,\sigma _1,\sigma _2=0}{\overset{q1}{}}}\gamma _{\sigma _1\sigma _2}^{\alpha \beta }\phi _{\alpha \sigma }^{}\phi _{\beta ,\sigma +\sigma _1}^{}\phi _{\alpha ,\sigma +\sigma _2}^{}\phi _{\beta ,\sigma +\sigma _1\sigma _2}^{},`$
for odd and even $`q`$, respectively, where $`\gamma _{mn}`$ and $`\gamma _{\sigma \sigma ^{}}^{\alpha \beta }`$ are arbitrary at this stage.
Imposing additional restrictions on the coefficients in the above Lagrangians from invariance under rotations and reflections and taking into account hermiticity and permutation symmetries, one may obtain a set of conditions on the $`\gamma `$ coefficients required to preserve the full triangular lattice symmetry. These conditions are given in Eqs. (A,A). For specific $`f`$, they can readily be solved to derive explicit forms for the Lagrangian. We will give these explicitly for $`f=1/2`$ and $`1/3`$ below.
Interestingly, it is possible to make at least one general observation concerning the symmetry properties of $`_1^{\mathrm{even}}`$. All terms in the continuum Lagrangian possess a global vortex $`U(1)`$ symmetry,
$$\phi _{\alpha \sigma }\phi _{\alpha \sigma }e^{i\theta }.$$
(29)
That is just a consequence of gauge invariance, expressing the conservation of the bosonic current. However, it is clear from Eq. (28) that the quartic Lagrangian $`_1^{\mathrm{even}}`$ possesses (at least) another, “staggered” U(1) symmetry:
$$\phi _{\alpha \sigma }\phi _{\alpha \sigma }e^{i\alpha \theta }.$$
(30)
This emergent $`U(1)`$ symmetry of the vortex theory implies that there are (at least) two conserved dual charges, that can be labelled by the index $`\alpha `$. As discussed in Ref. Balents04, , this is linked to the appearance of fractionally-charged bosonic excitations at the critical point. We will elaborate on the nature of these excitations in Sec. IV.
## III (Dual) Mean field theory
In this section we will discuss a mean field analysis of the vortex theory, focusing on the nature of the ordered Mott insulating states that occur. We first present some general aspects of how spatial order parameters are constructed in the vortex formalism, and give some physical picture of how to think of the different Mott states. The remaining two subsections describe the specific phase diagram in the case $`f=1/3`$ and the much more complicated and more interesting case $`f=1/2`$.
### III.1 Order parameters and Mott states
General argumentsOshikawa ; Hastings and physical reasoning seem to imply that, barring exotic situations such as phases with “topological order”,Wen Mott insulating states occuring for non-integral $`f`$ must break space group symmetries (and in particular translational symmetry). Such space group symmetry breaking is measured by spatial order parameters, the simplest of which (sufficient for our purposes) describe non-uniformity of the boson density (beyond that which is imposed by the underlying triangular substrate).
To visualize ordering patterns in the insulating phases we will find it convenient to introduce a general “density” function $`\rho (𝐫)`$, where $`𝐫`$ is a continuous real-space coordinate with $`𝐫=0`$ taken to coincide with a honeycomb lattice site of type “1” (Fig. 1). We construct $`\rho (𝐫)`$ to have the property that it transforms like a scalar boson density under all symmetry operations. It will be convenient to plot $`\rho (𝐫)`$ to graphically illustrate the symmetry of the non-uniform states that emerge in the theory. Writing $`𝐫=r_1𝐚_1+r_2𝐚_2`$, one can actually construct such a function quite generally for odd values of $`q`$:
$$\varrho (𝐫)=\underset{m,n}{}\varrho _{mn}\omega ^{mr_1+nr_2},$$
(31)
where the Fourier components $`\varrho _{mn}`$ serve as order parameters for different ordered states, and are given by:
$$\varrho _{mn}=S(m,n)\omega ^{mn/2+(nm)/6}\underset{\mathrm{}=0}{\overset{q1}{}}\omega ^m\mathrm{}\phi _{\mathrm{}}^{}\phi _{\mathrm{}+n}.$$
(32)
$`S(m,n)`$ here is a scalar form factor, that can not be determined from symmetry considerations, but should be chosen to depend only upon the magnitude of the wavevector $`m𝐛_1+n𝐛_2`$. A convenient and simple choice is the Lorentzian,
$$S(m,n)=\frac{1}{1+m^2+(m+2n)^2/3},$$
(33)
which we use only for plotting purposes. It is easy to check that $`\varrho _{mn}`$ indeed transform like Fourier components of density:
$`T_1:\varrho _{mn}\omega ^m\varrho _{mn},`$
$`T_2:\varrho _{mn}\omega ^n\varrho _{mn},`$
$`R_{2\pi /3}:\varrho _{mn}\varrho _{n,mn},`$
$`I_{d_1}:\varrho _{mn}\varrho _{n,m},`$
$`I_{d_2}:\varrho _{mn}\omega ^{(nm)/3}\varrho _{nm}.`$ (34)
A similar function can be constructed for $`q`$ even. When $`f=1/q`$, it takes the form
$$\varrho (𝐫)=\underset{m,n}{}\left[\varrho _{mn}^\alpha +\stackrel{~}{\varrho }_{mn}^\alpha e^{2\pi i\alpha (r_1+r_2)/3q}\right]\omega ^{mr_1+nr_2},$$
(35)
where the density wave amplitudes are given by
$`\varrho _{mn}^\alpha =S(m,n)e^{\pi i\alpha (nm)/3q}\omega ^{mn/2+(nm)/6}`$ (36)
$`\times `$ $`{\displaystyle \underset{\mathrm{}=0}{\overset{q1}{}}}\omega ^m\mathrm{}\phi _{\mathrm{}}^\alpha \phi _{\mathrm{}+n}^\alpha ,`$
$`\stackrel{~}{\varrho }_{mn}^\alpha =\stackrel{~}{S}(m,n)e^{i[\eta _2(\alpha )\eta _2(\alpha )]/2}\omega ^{mn/2+(nm)/6}`$
$`\times `$ $`{\displaystyle \underset{\mathrm{}=0}{\overset{q1}{}}}\omega ^m\mathrm{}\phi _{\mathrm{}}^\alpha \phi _{\mathrm{}+n}^\alpha .`$
Here the amplitude $`\stackrel{~}{S}(m,n)`$ has been taken in a simple form consistent with rotational symmetry:
$$\stackrel{~}{S}(m,n)=\frac{1}{1+(m+\alpha /3)^2+(m+2n+\alpha )^2/3}.$$
(37)
Eqs. (36) are written in terms of $`\eta _2(\alpha )`$, which also enters the PSG in general for even $`q`$, see Eqs. (A,107) of Appendix A. It is hard to determine for general $`q`$. However, for the case we will focus upon, $`f=1/2`$, one has
$$\eta _2(\alpha )=\frac{\pi \alpha }{12}\text{for }f=1/2.$$
(38)
We will present plots of $`\rho (𝐫)`$ for various mean field (and beyond, in the following section!) Mott insulating states.
These images, and their deconstruction into the Fourier amplitudes ($`\varrho _{mn}`$ etc.), characterize the broken spatial symmetry of the Mott states. They do not, however, directly give a physical picture for the ground state itself. Of course, for a general interacting boson model (and certainly in our approach where we do not specify the microscopic Hamiltonian), we cannot hope to write down the ground state wavefunction. Moreover, this has far too much information. What is conceptually useful is to understand how to write down a simple wavefunction appropriate to a Mott insulator with the same symmetries as predicted by the phenomenological theory. By an appropriate wavefunction, we mean one explicitly with the correct boson filling, with no long-range correlations beyond that of the Mott state, and consistent with incompressibility. We consider a satisfactory form to be a product state,
$$|\mathrm{\Psi }=\underset{\mathrm{}}{}|\psi (\mathrm{}).$$
(39)
The meaning of Eq. (39) is as follows. We divide the lattice into a set of non-overlapping identical clusters of sites – unit cells of the Mott state – labeled by the $`\mathrm{}`$. The $``$ indicates a direct product over states defined on the Hilbert space of each box, with the same state $`|\psi (\mathrm{})`$ chosen on each box. The state $`|\psi (\mathrm{})`$ must be an eigenstate of the total boson number on the given cluster, for this wavefunction to represent a Mott insulator, and this number should be chosen to match the filling, i.e. equal to $`f`$ times the number of sites in the cluster. Clearly this state has only local (within a cluster) charge fluctuations, consistent with incompressibility.
Of course in most cases many such wavefunctions can be constructed for a Mott state of a given symmetry, and moreover the true ground state wavefunction for a realistic hamiltonian will not have the direct product form. However, the above general construction can serve the purpose of providing an example of a state with the same symmetry properties as predicted by the phenomenological theory, and is expected to be in a sense (which we do not attempt to define precisely) adiabatically connected to all ground states in the Mott phase. To our knowledge, all well-established examples of bose Mott insulating ground states in the theoretical literature (e.g. on the square lattice, checkerboard, stripe, and VBS states) have such a construction. We therefore view the existence of such a wavefunction as a sort of consistency check on our results, and give examples in the following subsections. When the meaning is obvious, we give only a brief physical description of the state, with the understanding that one should keep the block product form, Eq. (39), in mind.
### III.2 $`f=1/3`$
At $`1/3`$-filling, the Mott state is not “frustrated” in the colloquial sense. The analysis below reveals that this case is extremely analogous to the superfluid-Mott transition on the square lattice at half-filling, and in particular provides another example of a possible deconfined quantum critical point of the type discussed in Ref. dcprefs, .
Solving Eq.(A) at $`f=1/3`$, we obtain the following form of the quartic potential in the continuum Lagrangian density:
$`_1=u\left(|\phi _0|^2+|\phi _1|^2+|\phi _2|^2\right)^2`$ (40)
$`+`$ $`v\{|\phi _0|^2|\phi _1|^2+|\phi _1|^2|\phi _2|^2+|\phi _2|^2|\phi _0|^2+2\text{R}e[e^{2\pi i/3}`$
$`\times `$ $`(\phi _0^2\phi _1\phi _2+\phi _2^2\phi _0\phi _1+\phi _1^2\phi _2\phi _0+c.c.)]\}.`$
As discussed in Ref. Balents04, , it is convenient to transform to a different set of variables that realize a “permutative representation” of the PSG. As in the square lattice case, such representations exist only for special fillings. In particular, it is easy to show that a permutative representation does not exist at $`f=1/2`$, in contrast to the square lattice case. At $`f=1/3`$ however, it does exist and is given by:
$`\xi _0={\displaystyle \frac{1}{\sqrt{3}}}\left(\phi _0+e^{2\pi i/3}\phi _1+\phi _2\right),`$
$`\xi _1={\displaystyle \frac{1}{\sqrt{3}}}\left(e^{2\pi i/3}\phi _0+\phi _1+\phi _2\right),`$
$`\xi _2={\displaystyle \frac{1}{\sqrt{3}}}\left(\phi _0+\phi _1+e^{2\pi i/3}\phi _2\right).`$ (41)
The physical meaning of the three vortex quantum numbers in the permutative representation of the PSG is that they represent three conserved (as shown below) dual charges (vorticities). The quantum numbers that are dual to these three vorticities are three real fractional (1/3 of the boson charge) charges. That is, a dual “vortex” in which the phase of any one of the $`\xi _{\mathrm{}}`$ fields winds by $`\pm 2\pi `$ at infinity carries a localized charge (boson number) of $`\pm 1/3`$ (see Ref. Balents04, for a simple derivation of this result in a more general context).
The representation of the PSG realized by the $`\xi _{\mathrm{}}`$ fields is permutative in that each symmetry operation is realized as the composition of a permutation of the $`\xi _{\mathrm{}}`$ fields and a simple phase rotation. In particular,
$`T_1:\xi _{\mathrm{}}\xi _{\mathrm{}+1},`$
$`T_2:\xi _{\mathrm{}}e^{2\pi (\mathrm{}+1)i/3}\xi _{\mathrm{}+1},`$
$`R_{2\pi /3}:\xi _0e^{\pi i/6}\xi _2,\xi _1i\xi _0,\xi _2e^{\pi i/6}\xi _1,`$
$`I_{d_1}:\xi _0e^{\pi i/6}\xi _2^{},\xi _1e^{\pi i/6}\xi _1^{},\xi _2e^{\pi i/6}\xi _0^{},`$
$`I_{d_2}:\xi _0e^{\pi i/6}\xi _0^{},\xi _1e^{\pi i/6}\xi _1^{},\xi _2e^{\pi i/6}\xi _2^{}.`$
The quartic potential simplifies greatly in these variables:
$`_1=u\left(|\xi _0|^2+|\xi _1|^2+|\xi _2|^2\right)^2`$ (43)
$`+`$ $`v\left(|\xi _0|^2|\xi _1|^2+|\xi _1|^2|\xi _2|^2+|\xi _2|^2|\xi _0|^2\right)`$
At the quartic level, it is immediately apparent that the microscopic overall $`U(1)`$ gauge symmetry required by the vortex non-locality has been elevated to a $`U(1)^3`$ symmetry under independent rotations of all three $`\xi _{\mathrm{}}`$ fields. Of this, only the group of equal rotations of all fields is gauge, leaving an additional $`U(1)^3/U(1)=U(1)\times U(1)`$ global (not gauge) symmetry of the dual theory. This emergent symmetry is broken at $`6^{\mathrm{th}}`$ order by the term
$$_2=w[(\xi _0^{}\xi _1)^3+(\xi _1^{}\xi _2)^3+(\xi _2^{}\xi _0)^3+c.c.].$$
(44)
The structure of the theory is thus rather similar to the vortex theory at $`f=1/2`$ on the square lattice.Lannert01 ; Balents04 It provides another example of deconfined criticality, if, as seems likely, the mean-field irrelevance of the higher-order term in $`_2`$ remains valid with fluctuations in $`2+1`$ dimensions, for some sign of $`v`$. The mean field phase diagram of the vortex theory can now be easily found analytically.
Let us now proceed with the mean field theory for $`f=1/3`$. For $`s<0`$ and $`|w|<|v|`$ one obtains 3 distinct insulating phases:
I. $`v\mathbf{>}\mathrm{𝟎}`$:
The energy is minimized if only one of the 3 vortex flavors condenses. This state is then clearly 3-fold degenerate and breaks all symmetries except reflection with respect to $`d_2`$ axis. In Fig.2 this state is visualized explicitly by plotting the corresponding density wave order parameter. This state is the simplest CDW state at $`1/3`$-filling, with an example wavefunction consisting of a boson number eigenstate on each site. Microscopically this phase is the natural Mott insulating state in a model in which the Mott transition is driven by strong nearest-neighbor repulsive interactions, like the XXZ model. We do not expect the superfluid-Mott transition to this state is likely to be continuous or deconfined when fluctuations are taken into account in $`2+1`$ dimensions.
II. $`v\mathbf{<}\mathrm{𝟎}`$:
It is energetically favorable to condense all 3 vortex flavors, so that all vortex fields have equal magnitude. There are then 2 different phases, depending on the sign of $`w`$. These states are those expected to be connected by a deconfined quantum critical point to the superfluid state.
1. $`w<0`$.
Writing $`\xi _{\mathrm{}}|\xi |e^{i\theta _{\mathrm{}}}`$, the minimum is achieved when:
$$\theta _1\theta _0=2\pi m/3,\theta _2\theta _0=2\pi n/3,$$
(45)
for $`m,n=0,1,2`$.
This state is thus 9-fold degenerate and corresponds to period-3 site-centered stripes, running in $`𝐚_1,𝐚_2`$ and $`𝐚_1+𝐚_2`$ directions, see Fig.3. One may construct a wavefunction for this state by e.g. taking linear 3-site clusters at a $`60^{}`$ angle to the stripe (e.g. horizontal in Fig. 3), and putting a boson at the center of the cluster, leaving the other sites empty.
2. $`w>0`$.
One can readily verify that the ground state is 18-fold degenerate, the distinct solutions being obtained from
$$\theta _1\theta _0=\theta _2\theta _0=2\pi /9,$$
(46)
by applying translations and the reflection $`I_{d_2}`$. The corresponding characteristic density pattern is shown in Fig.4. It is adiabatically connected to a “bubble” phase or crystalline state in which one boson is placed on each site of each elementary triangle (3 bosons total) on a $`3\times 3`$ triangular superlattice. The 18-fold degeneracy results from 9 states obtained by translating this pattern, and another 9 states obtained by choosing say down-pointing instead of up-pointing triangles.
### III.3 $`f=1/2`$
#### III.3.1 Action and symmetries
We now turn to the more complicated and interesting case of $`f=1/2`$. It is straightforward to show (by a simple generalization of the argument used in Ref. Balents04, to prove the absence of a permutative representation in the $`f=1/3`$ case on the square lattice), that for this case there is no permutative representation of the PSG. Unlike in the square lattice case, there are also, surprisingly, more low-energy vortex modes for $`f=1/2`$ case than for $`f=1/3`$. The emergent low-energy symmetry amongst these modes is, moreover, nonabelian, with an $`SU(2)`$ “pseudo-spin” subgroup that we will uncover below. This structure has interesting consequences for the phases and excitations that will be explored here and in the following section in some detail.
The quartic potential can be found by solving Eq.(A). It turns out to have the following form:
$`_1=u(\phi _{\alpha \sigma }^{}\phi _{\alpha \sigma }^{})^2+v(\phi _{+\sigma }^{}\phi _{+\sigma }^{})(\phi _\sigma ^{}^{}\phi _\sigma ^{}^{})`$ (47)
$`+w(|\phi _0|^2|\phi _{+0}|^2+\phi _1^{}\phi _{+1}^{}\phi _0\phi _{+0}+(01)).`$
To make the symmetries of the Lagrangian more transparent, it is useful to rewrite it as follows. First we define $`z_{\alpha \sigma }`$ pseudospinor variables,
$`\phi _\sigma `$ $`=`$ $`ϵ_{\sigma \sigma ^{}}z_\sigma ^{},`$
$`\phi _{+\sigma }`$ $`=`$ $`z_\sigma ,`$ (48)
where $`ϵ_{\sigma \sigma ^{}}`$ is the antisymmetric tensor with $`ϵ_{01}=ϵ_{10}=1`$. The quartic action has $`SU(2)`$ symmetry under rotations of the $`\sigma `$ index of $`z_{\alpha \sigma }`$. This is made manifest by introducing the pseudospin vector
$$S_\alpha ^a=z_{\alpha \sigma }^{}\tau _{\sigma \sigma ^{}}^az_{\alpha \sigma ^{}}^{},$$
(49)
where $`\tau ^a,a=x,y,z`$ are the Pauli matrices, and summation on the repeated $`\sigma ,\sigma ^{}`$ indices is implied. The transformation properties of the $`𝐒_\pm `$ vectors are particularly simple, and given in Appendix B. In terms of these pseudospin variables the quartic potential becomes, after a trivial redefinition of $`v`$ and $`w`$ couplings:
$$_1=u\left(S_++S_{}\right)^2+vS_+S_{}+w_1𝐒_+𝐒_{},$$
(50)
where $`S_\alpha =|𝐒_\alpha |=\sqrt{z_{\alpha \sigma }^{}z_{\alpha \sigma }^{}}`$ (sum on $`\sigma `$ implied). It is now clear that the quartic potential has, in addition to the microscopic gauge $`U(1)`$ symmetry, an $`SU(2)\times U(1)\times Z_2`$ invariance. The $`SU(2)`$ symmetry is manifest in Eq.(50), the extra $`U(1)`$ symmetry is the “staggered” $`U(1)`$, already mentioned above, see Eq.(30), also manifest since $`S_\alpha ^a`$ are independent of the staggered $`U(1)`$ phase. The $`Z_2`$ symmetry is the interchange $`𝐒_+𝐒_{}`$ (particle-hole symmetry $`C`$ in fact requires this invariance up to a sign, though Eq. (50) is obtained without using $`C`$). The pseudospin variables are directly related to the $`\varrho _{mn}^\alpha `$ density components:
$`\varrho _{00}^\alpha =S_\alpha ,`$
$`\varrho _{01}^\alpha =\alpha e^{\pi i(1+\alpha )/6}S_\alpha ^x,`$
$`\varrho _{10}^\alpha =\alpha e^{\pi i(1+\alpha )/6}S_\alpha ^z,`$
$`\varrho _{11}^\alpha =S_\alpha ^y.`$ (51)
The physical import of the $`SU(2)`$ symmetry of Eq. (50) is now clear: all CDW states, related to each other by arbitrary rotations in the space of the three Fourier components of the density, are degenerate at this order. This degeneracy will be weakly broken, of course, by higher order terms in the action. As discussed in the introduction, the large emergent symmetry is thereby connected with geometrical charge frustration at $`f=1/2`$.
#### III.3.2 Order parameters
The pseudospin vectors $`𝐒_\pm `$ serve as gauge-invariant order parameters to characterize the breaking of the $`SU(2)`$ symmetry. It is instructive to construct two other such order parameters. The emergent $`U(1)`$ symmetry is best characterized by a complex order parameter, $`\psi `$, defined by
$$\psi =e^{i\pi /4}z_{+\sigma }^{}z_\sigma ^{},$$
(52)
where we have included the $`\pi /4`$ phase factor for later convenience. One may also define an Ising order paramer $`\mathrm{\Phi }`$, with
$$\mathrm{\Phi }=|z_+|^2|z_{}|^2.$$
(53)
Non-zero $`\mathrm{\Phi }`$ implies $`I_{d_2}`$ and $`C`$ are broken. The physical meaning of non-vanishing $`𝐒_\pm `$ and $`\psi `$ will be elucidated in detail in the following.
Though the terms which break these symmetries are small near the Mott QCP (hopefully irrelevant there), they are important at sufficiently low energy. We must therefore consider those higher order terms in the action which are required to reduce the $`SU(2)\times U(1)`$ symmetry to only what is required by the PSG. Higher order terms which do not reduce this symmetry need not be considered.
First consider the $`SU(2)`$ pseudospin symmetry. In the absence of microscopic particle-hole invariance, it is broken at $`6^{\mathrm{th}}`$ order by a term of the form $`S_+^xS_+^yS_+^z+S_{}^xS_{}^yS_{}^z`$. We will, however, for concreteness focus on the case relevant to the XXZ model and other pairwise interacting boson lattice models, in which particle-hole symmetry is an invariance of the theory. In this case, the $`SU(2)`$ symmetry is broken only at the 8th order by a “cubic anisotropy” term,
$$_2=w_2\underset{\alpha }{}\left[\left(S_\alpha ^x\right)^4+\left(S_\alpha ^y\right)^4+\left(S_\alpha ^z\right)^4\right].$$
(54)
The “staggered” $`U(1)`$ symmetry is more persistent and is broken at the 12th order. The simplest term at this order that breaks the “staggered” $`U(1)`$ symmetry is:
$`_3=w_3\mathrm{Re}\left(\psi ^6\right).`$ (55)
It is useful to reorganize various terms in the density expansion of Eq. (35) to understand in more detail the nature of the different order parameters. The wavevectors referred to in the following are labelled in Fig. 5.
Consider first Ising order. Non-zero $`\mathrm{\Phi }`$ implies only that the reflection $`I_{d_2}`$ and particle-hole symmetry $`C`$ are broken. Thus it corresponds only to a modulation of the density within the primitive unit cell of the triangular lattice. The corresponding density modulation therefore occurs entirely at reciprocal lattice vectors. The smallest set of reciprocal lattice wavevectors that can describe the modulation are $`𝐁_0=(0,2\pi )`$, $`𝐁_0=(2\pi ,2\pi )`$, $`𝐁_2=(2\pi ,0)`$. This density modulation is
$$\rho _\mathrm{\Phi }(𝐫)=\mathrm{\Phi }\underset{n=0,1,2}{}\mathrm{cos}(𝐁_n𝐫+\frac{5\pi }{6}).$$
(56)
While $`\rho _\mathrm{\Phi }=0`$ on triangular (direct) lattice sites (as it must), it alternates sign on sites of the dual honeycomb lattice, i.e. takes opposite signs on centers of triangles of the direct lattice.
Now let us turn to pseudospin ordering. The existence of non-vanishing $`𝐒_\pm `$ implies charge ordering at the wavevectors $`𝐆_x=(0,\pi )`$, $`𝐆_y=(\pi ,\pi )`$, $`𝐆_z=(\pi ,0)`$, which lie at the centers of the zone edges. In particular, the associated density modulations take the form
$$\rho _S(𝐫)=\underset{a=x,y,z}{}\mathrm{Re}\left[(S_+^ae^{i\pi /3}S_{}^a)e^{i𝐆_a𝐫}\right],$$
(57)
neglecting higher harmonics which do not change the symmetry of $`\rho _S(𝐫)`$.
Next consider the XY order parameter $`\psi `$. A state with $`\psi `$ exhibits a three-sublattice structure, characterized by the zone boundary wavevectors $`𝐐_0=2\pi (\frac{2}{3},\frac{1}{3})`$, $`𝐐_1=2\pi (\frac{1}{3},\frac{2}{3})`$, $`𝐐_2=2\pi (\frac{1}{3},\frac{1}{3})`$:
$$\rho _\psi (𝐫)=\mathrm{Re}\left[\psi e^{i\pi /6}\underset{n=0,1,2}{}e^{i2\pi n/3}e^{i𝐐_n𝐫}\right].$$
(58)
In fact, these three wavevectors differ only by reciprocal lattice vectors. From this, it is straightforward to show that $`\rho _\psi (𝐫)`$ vanishes on dual lattice sites, so that all triangular plaquettes of the direct lattice are equivalent up to rotations in an XY ordered state.
Finally, simultaneous breaking of pseudospin and XY symmetry is characterized by the “composite” order parameter $`𝐝`$, a complex vector, defined as
$$𝐝=z_{+\sigma }^{}𝝉_{\sigma \sigma ^{}}z_\sigma ^{}^{}.$$
(59)
When $`𝐝0`$, density modulations appear at the wavevectors $`𝐊_x=2\pi (\frac{1}{3},\frac{1}{6})`$, $`𝐊_y=2\pi (\frac{1}{6},\frac{1}{6})`$, $`𝐊_z=2\pi (\frac{1}{6},\frac{1}{3})`$, which lie within the zone. The corresponding density is
$$\rho _d(𝐫)=\mathrm{Re}\left[e^{3\pi i/4}\underset{a=0,1,2}{}d_ae^{2\pi (a1)i/3}e^{i𝐊_a𝐫}\right],$$
(60)
where we have identified $`a=x,y,z`$ with $`a=0,1,2`$ respectively.
#### III.3.3 Mean field phases
The mean field phase diagram of $`_0+_1+_2+_3`$ can be easily obtained analytically. One finds 10 different phases: 2 for $`v>0`$ and 8 for $`v<0`$. We will not attempt to be exhaustive in describing these states. We will, however, discuss some of the phases in detail, and go into some general aspects of the 8 cases with $`v<0`$ in Sec. IV. Minimizing the mean-field energy functional, $`v<0`$ implies that both $`\alpha =\pm `$ vortex pseudospinors are condensed with equal amplitude, so $`S_+=S_{}`$. The sign of $`w_1`$ determines the relative pseudospin orientation. For $`w_1<0`$, they are parallel, i.e. $`𝐒_+=𝐒_{}𝐒`$. In terms of the vortex variables this condition is most generally solved by
$$z_{\pm \sigma }=z_\sigma e^{\pm i\theta /2},$$
(61)
where $`𝐒=z_\sigma ^{}𝝉_{\sigma \sigma ^{}}z_\sigma ^{}^{}`$. For $`w_1>0`$, the two pseudospin vectors are antiparallel, $`𝐒_+=𝐒_{}𝐒`$, which implies
$`z_{+\sigma }`$ $`=`$ $`z_\sigma e^{i\theta /2},`$
$`z_\sigma `$ $`=`$ $`ϵ_{\sigma \sigma ^{}}^{}z_\sigma ^{}^{}e^{i\theta /2}.`$ (62)
We note that the dual gauge symmetry acts differently in the two cases. For parallel pseudospins, a gauge transformation takes $`z_\sigma e^{i\chi }z_\sigma `$, while for antiparallel pseudospins, instead $`\theta \theta +2\chi `$.
The remaining terms, $`w_2`$ and $`w_3`$, fix the remaining non-gauge symmetries. The sign of $`w_2`$ chooses the easy axes of the pseudospin vector, along $`(100)`$ and symmetry-related axes for $`w_2<0`$, and along $`(111)`$ and related axes for $`w_2>0`$. The above conditions leave only the relative phase between the spinors $`z_\pm `$, corresponding to the “staggered” $`U(1)`$ symmetry of the corresponding terms in the vortex Lagrangian.
This relative phase is fixed at $`12^{\mathrm{th}}`$ order in $`z_{\pm \sigma }`$, for instance for $`w_1<0`$ by the $`_3`$ term (for $`w_1>0`$, a more complicated $`12^{\mathrm{th}}`$ order term must be included as the $`w_3`$ interaction vanishes in that case). Depending on the sign of $`w_3`$, the energy minimum is achieved when $`\mathrm{sin}(6\theta )=\pm 1`$, i.e.
$$\theta =\frac{\pi }{12}+\frac{\pi n}{3},n=0,\mathrm{},5,$$
(63)
or
$$\theta =\frac{\pi }{4}+\frac{\pi n}{3},n=0,\mathrm{},5,$$
(64)
One of the states, corresponding to $`\theta =\pi /12`$ and
$`z_{+0}`$ $`=`$ $`z_{+1}={\displaystyle \frac{e^{i\theta /2}}{\sqrt{2}}},`$
$`z_0`$ $`=`$ $`z_1={\displaystyle \frac{e^{i\theta /2}}{\sqrt{2}}},`$ (65)
is shown in Fig.6. We will elucidate the physics of this particular pair of states in some detail in the discussion section.
We can generally classify all the mean field states by their degeneracies and the corresponding unit cell sizes. The states with the pseudospin easy axis along (100) and parallel pseudospins (the state in Fig.6 belongs to this group) are all 36-fold degenerate and have a 6-site unit cell. States with parallel pseudospins but with the easy axis along (111) are 48-fold degenerate and have the largest, 12-site unit cell. In the case of antiparallel pseudospins, ground state degeneracies are the same, but the unit cell size of the 36-fold degenerate states doubles to 12 sites. States with the smallest unit cells are obtained when only one of the pseudospinors is condensed. In this case one obtains a 6-fold degenerate ground state and a 2-site unit cell when (100) is the easy axis, and an 8-fold degenerate state with a 4-site unit cell when the easy axis is along the (111) direction.
For the phases of most interest ($`v<0`$) in which the pseudospin vectors are either parallel or antiparallel, some of the density functions associated to the order parameters in Sec. III.3.2 can be simplified to an extent. The only qualitative case is the pseudospin vector density $`\rho _S`$. When $`𝐒_+=𝐒_{}=𝐒`$, it reduces to
$$\rho _S(𝐫)=\underset{a=x,y,z}{}S^a\mathrm{cos}(𝐆_a𝐫\frac{2\pi }{3}).$$
(66)
Most interesting, when $`𝐒_+=𝐒_{}=𝐒`$, one has instead
$$\rho _{S,\parallel ̸}(𝐫)=\underset{a=x,y,z}{}S^a\mathrm{cos}(𝐆_a𝐫\frac{\pi }{6}).$$
(67)
In this case, it is noteworthy that $`\rho _{S,\parallel ̸}`$ vanishes on all triangular lattice sites. This is a consequence of the fact that a configuration of antiparallel pseudospins preserves particle-hole symmetry, so modulations can occur only in bond or plaquette “kinetic” terms. Thus $`𝐒_+𝐒_{}`$ may be considered a (particular) purely valence bond solid order parameter.
The composite order parameter $`𝐝`$ also simplifies once the pseudospin order is determined. For parallel pseudospins, one simply has $`𝐝=𝐒`$. For antiparallel pseudospins, instead, one has $`𝐝\times 𝐝^{}=2i𝐒`$, which implies
$`𝐝=𝐞_1i𝐞_2,𝐞_1\times 𝐞_2=𝐒,`$ (68)
so that $`𝐞_1,𝐞_2,𝐒`$ form a right-handed orthogonal frame in the O(3) spin space. The angle of $`𝐞_1`$ in the XY plane is arbitrary, and determined by (twice) the phase of $`z_\sigma `$.
## IV “Hard Spin” Description: Beyond Mean-Field Theory
In the preceding section, we have followed a Landau-theory like procedure (albeit with non-LGW vortex fields) in expanding the effective action in a power series in the $`\phi _{\mathrm{}}`$ fields, whose amplitude is viewed as small in the vicinity of the Mott QCP. In low-dimensional statistical mechanics, it is often preferable to formulate the theory in terms of “hard spin” variables, in which the amplitude of the order parameter field(s) is fixed, and only the “angular” degrees of freedom are free to fluctuate and vary in space and time. Examples include the Kosterlitz-Thouless theory of the XY phase transition, and the non-linear sigma model formulation of $`O(n)`$ models. The intuitive rationale for such an approach is that, in low dimensions, fluctuations suppress the ordering point of the transition well below the mean-field point, so that substantial amplitude is already developed in the true critical region. Whatever the rationale, there are some advantages to such a hard-spin approach. Duality transformations generally apply to hard-spin models. Hard-spin variables are particularly appropriate to describe the elementary excitations of “ordered” phases in which the amplitude of the fields is on average large, and only the Goldstone-like fluctuations of the orientation of these fields comprise low-energy excitations. Finally, a hard-spin formulation returned to the lattice is fully-regularized, and can thereby address non-perturbative phenomena in a controlled manner.
In this section, we will provide and analyze a hard-spin formulation of the dual vortex action. These allow us to identify the excitations and their quantum numbers within the Mott phases. Notably, we find that the most interesting Mott states support two distinct kinds of excitations. First, there are $`1/2`$-charged (“spinon” in the spin-$`1/2`$ XXZ language) “vortex” excitations which are linearly confined in pairs deep in the Mott state, but are only logarithmically interacting up to a long “confinement length” near the superfluid-Mott transition. Second, there are additional unit charged “skyrmion” excitations which are everywhere deconfined in the Mott state. These are adiabatically connected to single boson vacancies/interstitials in the Mott solid, but become topological as the Mott transition is approached. The hard-spin models also provide a firm ground on which to study other phases, notably supersolids, which do not occur within a mean field treatment of the vortex field theory, but are extremely natural in this approach.
### IV.1 Formulation of hard spin model
To write down an appropriate hard-spin model, we imagine tuning $`s<0`$ to a point beyond the mean field Mott transition point. At such a point, the minimum action configurations have non-zero amplitude. We will assume the mean amplitude is determined by a balance between the quadratic Lagrangian, Eq. (25) and the quartic terms in $`_1`$, Eq. (50). The minimum action saddle points of the combination of these two terms are constant in space-time, but allow for a continuous set of orientations in the field space. We focus in particular on $`v<0`$, so that the magnitude $`S_+=S_{}`$ is fixed. We will focus primarily on the case $`w_1<0`$, so that the saddle point has $`𝐒_+=𝐒_{}`$. This is the most interesting case, because, as we shall see, the recently-determined supersolid phase of the XXZ model can be understood in this framework. At the end of this section, we will briefly summarize the results of similar analysis for $`w_1>0`$, corresponding to anti-parallel pseudospins.
We further suppose that $`s`$ is negative enough that fluctuations in the above conditions may be neglected, but that within these constraints the $`z_{\alpha \sigma }`$ fields can vary spatially. We will therefore absorb any effects of the magnitude of the fields into coefficients, and without loss of generality normalize to $`S_+=S_{}=1`$. Ultimately, the higher order terms will still be included, but can be considered small perturbations.
The most general solution of $`𝐒_+=𝐒_{}=𝐒`$ constraint in terms of the vortex variables is given by Eq. (61). We will therefore rewrite the action in terms of the CP<sup>1</sup> field $`z_\sigma `$ and an XY field $`e^{i\theta /2}`$. Note that this solution possesses a $`Z_2`$ gauge invariance under
$$z_\sigma z_\sigma ,\theta \theta +2\pi .$$
(69)
This is in addition to the physical symmetries of the model. It is a gauge invariance since it can be performed independently at each space-time point without changing $`z_{\alpha \sigma }`$, and hence physical quantities.
Inserting Eq. (61) into the action, and regularizing it on a space-time lattice, we obtain
$$_{\mathrm{HS}}^{}=t_ve^{iA_{i\mu }}z_{i\sigma }^{}z_{i+\mu \sigma }\mathrm{cos}(\mathrm{\Delta }_\mu \theta _i/2)+\frac{1}{2e^2}(ϵ_{\mu \nu \lambda }\mathrm{\Delta }_\nu A_{i\lambda })^2,$$
(70)
with $`z_{i\sigma }^{}z_{i\sigma }^{}=1`$ normalized on each site $`i`$ of the space-time lattice. Note that Eq. (70) is indeed invariant independently under Eq. (69) at each point $`i`$. It is convenient to rewrite the first term in Eq. (70), making the Ising gauge symmetry explicit by introducing an Ising gauge field $`\sigma _{i\mu }`$ which resides on the link $`(i,i+\mu )`$:
$`_{\mathrm{Z}_2}^{}`$ $`=`$ $`t_z\sigma _{i\mu }e^{iA_{i\mu }}z_{i\sigma }^{}z_{i+\mu \sigma }t_\theta \sigma _{i\mu }\mathrm{cos}(\mathrm{\Delta }_\mu \theta _i/2)`$ (71)
$`+{\displaystyle \frac{1}{2e^2}}(ϵ_{\mu \nu \lambda }\mathrm{\Delta }_\nu A_{i\lambda })^2,`$
where we have introduced two independent parameters $`t_z,t_\theta `$. One could add a plaquette interaction (line product of $`\sigma _{i\mu }`$ around plaquettes), which is a standard “kinetic term” for the $`Z_2`$ gauge field for further generality. We are, however, most interested in the limit in which it is absent. In this case, one may sum over $`\sigma _{i\mu }`$ on each link independently, to return to an action of the form of Eq. (70) (with additional higher-order terms). We will not attempt, however, to constrain Eq. (71) to be exactly equivalent to Eq. (70). Instead, since we are anyway constructing a phenomenological theory, we regard the freedom to vary $`t_z,t_\theta `$ independently as a means of capturing the different possible tendencies due to fluctuation effects and details of microscopic dynamics in different physical systems.
### IV.2 Phase diagram for parallel pseudospins
Let us now discuss the phase diagram of Eq. (71). For $`t_z,t_\theta 1`$, both $`\theta _i`$ and $`z_{i\sigma }`$ variables are disordered, and the dual gauge field $`A_{i\mu }`$ is gapless. This is the superfluid phase. For $`t_\theta t_z1`$ large and of the same order, we expect that both the CP<sup>1</sup> and XY variables are ordered. This is the Mott insulator, whose precise nature depends upon the anisotropy terms we have neglected to write.
Now suppose $`t_\theta 1`$ but $`t_z1`$. In this limit, we expect the XY variables condense. This is a “Higgs” phaseKogut from the point of view of the $`Z_2`$ gauge variables: the linear coupling to $`\mathrm{cos}(\mathrm{\Delta }_\mu \theta _i/2)`$ means that the $`\sigma _{i\mu }`$ fields can be regarded as having some non-vanishing expectation value in this state. The CP<sup>1</sup> fields however remain uncondensed, and $`A_{i\mu }`$ remains gapless, so this state retains superfluidity. It does, however, break spatial symmetry. To understand the nature of this symmetry breaking, let us take for simplicity $`t_z=0`$, and again imagine “summing out” the $`Z_2`$ gauge fields while varying $`t_\theta `$. The effective Lagrange density is then
$`_{XY}`$ $`=`$ $`V(\mathrm{\Delta }_\mu \theta _i)+{\displaystyle \frac{1}{2e^2}}(ϵ_{\mu \nu \lambda }\mathrm{\Delta }_\nu A_{i\lambda })^2`$ (72)
with $`V(\mathrm{\Theta })=\mathrm{ln}\mathrm{cosh}\left[\frac{t_\theta }{\sqrt{2}}\sqrt{1+\mathrm{cos}(\mathrm{\Theta })}\right]`$. This somewhat unconventional gradient term has all the same symmetries as the usual $`\mathrm{cos}(\mathrm{\Delta }_\mu \theta _i)`$ term in an XY model, and indeed reduces to that form for small $`t_\theta `$. On increasing $`t_\theta `$, therefore, a 3D=2+1 dimensional XY transition is expected, into a state with a non-zero expectation value of $`e^{i\theta }`$. Comparison with Eq. (52) indicates that $`\psi e^{i\pi /4}e^{i\theta }`$ is an order parameter for this transition. Note that $`e^{i\theta _i/2}`$ is not an order parameter since it is not ($`Z_2`$) gauge invariant. As noted earlier, $`\psi `$ is exactly the order parameter identified in Refs. Melko05, ; Damle05, ; Troyer05, as characterizing the supersolid phase of the XXZ model on the triangular lattice. This supersolid has a three-sublattice structure, so we will denote it by SS3. Actually there are two different ordering patterns possible within the tripled unit cell, depending upon the sign of the $`6`$-fold anisotropy term, $`w_3`$, which should be added to Eq. (72). They are shown in Fig. 7.
Finally, consider similarly the situation when $`t_\theta 1`$ but $`t_z`$ varies from small to large. For large $`t_z`$, we then expect the CP<sup>1</sup> variables order but the XY variables remain uncondensed. Analogously to the previous case, imagine increasing $`t_z`$ from small to large with $`t_\theta =0`$. One can again integrate out the Ising gauge field to obtain an effective action which is $`Z_2`$-gauge invariant. In this case, there are two distinct types of “kinetic” terms which arise on nearest-neighbor bonds. For small $`t_z`$, they take the form
$`_{SU(2)}`$ $`=`$ $`t_S𝐒_i𝐒_jt_2e^{2iA_{i\mu }}z_{i\sigma }^{}z_{i\sigma ^{}}^{}z_{i+\mu ,\sigma }^{}z_{i+\mu ,\sigma ^{}}^{}`$ (73)
$`+{\displaystyle \frac{1}{2e^2}}(ϵ_{\mu \nu \lambda }\mathrm{\Delta }_\nu A_{i\lambda })^2,`$
where $`t_S,t_2t_z^2`$ and $`𝐒_i=z_i^{}𝝉z_i^{}`$. Two distinct types of “orderings” are clearly possible on increasing $`t_z`$. The most natural possibility, driven by $`t_S`$, is for $`𝐒`$ to order. As above, this is a supersolid, but with a different set of possible charge order patterns, characterized by the zone boundary center wavevectors rather than those at zone corners. An alternative possibility, driven by $`t_2`$, is that the vortex pair field $`z_{i\sigma }z_{i\sigma ^{}}`$ condenses. Such a paired vortex condensate is a spin-liquid insulator (see Ref. Balents00, ) (a “$`Z_2`$” spin liquid in the now-conventional nomenclatureSenthilFisher ; Wen91 ). The $`A_{i\mu }`$ gauge fluctuations will tend to suppress such pair field condensation, so we expect the supersolid phase with $`𝐒0`$ to occur first on increasing $`t_z`$. Such supersolid states have a maximum period of $`2`$ lattice sites along the principle axes of the triangular lattice (as can be seen from the behavior of translations in Eqs. (142), so we denote these phases by SS2 (they may have doubled or quadrupled unit cells, depending upon the orientation of the pseudospin vector). The two different ordering patterns for different signs of $`w_2`$ are shown in Fig. 8.
Putting the different limits of this analysis together and making the simplest possible interpolation, we expect the phase diagram in Fig. 9. The Mott state may be reached from the superfluid in at least three distinct ways: by a direct transition described by the continuum vortex Lagrangian in the previous section, or via two distinct intermediate supersolid phases. The transitions from the superfluid to the two supersolids are “conventional”, i.e. of LGW type, since they are described by ordering of the gauge-invariant order parameters $`\psi `$ and $`𝐒`$. The transitions from the supersolids, by contrast, are unconventional. This is clear from the fact that the Mott insulator differs from either supersolid by breaking more spatial symmetry and by having no off-diagonal long range order, i.e. by being non-superfluid. Thus two symmetry-unrelated order parameters must change in these transitions. We will return to the nature of these transitions after first discussing the elementary excitations of the different phases.
### IV.3 Elementary excitations
#### IV.3.1 Superfluid phase
The hard spin model is convenient for describing the elementary excitations of the phases discussed above. First consider the superfluid. In this case, the elementary excitations are simply vortices, and the vortex field theory of the previous two sections already gives a description of the elementary vortex multiplet, consisting in this case of $`4`$ vortex flavors (carrying pseudospin-$`1/2`$ and XY “charge” $`\pm 1/2`$). This should be reproduced by the hard spin model. Naïvely, the “particles” of the hard spin model are created separately by the $`z_\sigma `$ and $`e^{\pm i\theta /2}`$ fields, so carry only one or the other of pseudospin or XY charge. However, in the superfluid region of Fig. 9, the $`Z_2`$ gauge charges (whose interactions are mediated by $`\sigma _{i\mu }`$) are confined, so the true elementary excitations are $`Z_2`$ gauge-neutral bound states $`e^{\pm i\theta /2}z_\sigma `$ which have precisely the appropriate quantum numbers of the $`z_{\pm \sigma }`$ vortices.
#### IV.3.2 Mott phase
Next consider the Mott state. In reality this comprises a number of different phases, depending upon the signs of the anisotropies $`w_2,w_3`$. However, near to the superfluid-Mott transition, the latter terms are small, and these distinct phases are approximately unified into one continuous manifold. It is useful to discuss the elementary excitations therefore in the same approximation, and afterward describe how they are modified once anisotropy is included. From the point of view of the $`Z_2`$ hard spin model, Eq. (71), the Mott state is a Higgs phase, with both XY and CP<sup>1</sup> fields condensed. The $`Z_2`$ gauge field can be regarded, in a choice of gauge, as uniform $`\sigma _{i\mu }1`$, in the ground state. If we neglect the possibility of deforming the $`Z_2`$ gauge field in excited states, there are then two “obvious” topological excitations – time independent solitons that behave as quantum particles – corresponding to textures in the CP<sup>1</sup> and XY fields.
First consider the CP<sup>1</sup> field. At spatial infinity, $`z_{i\sigma }`$ must vary slowly in space to maintain minimal action, as must $`A_{i\mu }`$ for the same reason. We may therefore take a continuum limit of the first term in Eq. (71), and, taking $`\sigma _{i\mu }1`$, one finds the Lagrangian
$`_z`$ $``$ $`{\displaystyle \frac{1}{2}}\varrho _s|_\mu 𝐒|^2+\kappa (𝒞_\mu A_\mu )^2`$ (74)
$`+{\displaystyle \frac{1}{2e^2}}(ϵ_{\mu \nu \lambda }_\nu A_{i\lambda })^2`$
with $`\varrho _s,\kappa t_z`$, and
$$𝒞_\mu =\frac{i}{2}\left(z_\sigma ^{}_\mu z_\sigma ^{}_\mu z_\sigma ^{}z_\sigma ^{}\right).$$
(75)
Minimum action configurations therefore have $`_\mu 𝐒=0`$ and $`A_\mu =𝒞_\mu `$ at infinity. This requires that, at infinity, only the phase of the spinor varies, i.e. $`z_\sigma =\xi _\sigma e^{i\mathrm{\Theta }}`$, where $`\xi _\sigma `$ is a constant normalized spinor, and $`\mathrm{\Theta }`$ may depend upon the polar angle from the origin. At infinity, then
$$𝒞_\mu =_\mu \mathrm{\Theta }.$$
(76)
Single-valuedness of $`z_\sigma `$ allows topologically non-trivial configurations in which $`\mathrm{\Theta }`$ winds by an integer multiple of $`2\pi `$, whence
$$𝑑x_\mu 𝒞_\mu =𝑑x_\mu A_\mu =2\pi n_s,$$
(77)
the integer $`n_s`$ being a topological index. Thus the dual flux of such configurations is quantized in units of the dual flux quantum. Since the dual flux measures physical charge, such solitons are particle-like excitations of the Mott state with physical integral boson charge $`n_s`$. While at infinity the spinor varies only through $`\mathrm{\Theta }`$, (and hence the pseudospin $`𝐒`$ is constant), this cannot hold everywhere in space, since such a “vortex” in $`\mathrm{\Theta }`$ must have a singularity somewhere. If the spinor is assumed to vary everywhere slowly in space, so that a uniform continuum limit can be taken everywhere, then the singularity is avoided by having the associated amplitude vanish at the skyrmion’s “core”, e.g. in a configuration of the form
$$z_\sigma =f(r)e^{i\mathrm{\Theta }}\xi _\sigma +\sqrt{1f(r)}\eta _\sigma ,$$
(78)
where $`\xi _\sigma `$ and $`\eta _\sigma `$ are two normalized orthogonal constant spinors, $`\xi _\sigma ^{}\eta _\sigma ^{}=0`$, and $`f(r)1`$ as $`r\mathrm{}`$, $`f(r)0`$ as $`r0`$ ($`r`$ is the radial coordinate from the skyrmion center). The non-collinear variation of $`z_\sigma `$ indicates a non-trivial texture of the pseudospin in the skyrmion. Quite generally, if $`z_\sigma `$ is analytic, one can show that
$$n_s=\frac{1}{4\pi }d^2r𝐒_x𝐒\times _y𝐒,$$
(79)
directly relating the skyrmion number to the pseudospin texture. Because the pseudospin itself is constant at infinity, it is apparent that the skyrmion has a finite size, in the example of Eq. (78) determined by the range of significant spatial variation of $`f(r)`$. This scale is dependent upon details of the Hamiltonian within the Mott phase, and in general the above description of the spatially-varying pseudospin is valid only if this scale is much larger than the unit cell of the Mott charge ordering pattern. Deep in the Mott phase (i.e far from the Mott transition), this may not be the case, and in that case there is not necessarily any sharp meaning to the pseudospin texture. In this sense, the skyrmion/antiskyrmion can be considered as adiabatically connected to a simple, and patently non-topological, vacancy or interstitial defect of the Mott “solid”.
Near to the Mott transition, however, the skyrmion is expected to be large, as we now show. The size of the skyrmions is determined by the balance of the anisotropy energy $`_2`$ and the interaction energy $`(ϵ_{\mu \nu \lambda }_\nu A_\lambda )^2`$. A rough estimate for the skyrmion size can be obtained by a simple dimensional analysis. The anisotropy energy of a skyrmion of size $`\lambda `$ is of the order of $`w_2|S|^4\lambda ^2`$, where $`|S|`$ is the unrescaled amplitude of the pseudospin vector order parameter. On the other hand, the interaction energy is of the order $`1/e^2\lambda ^2`$. The optimal skyrmion size is therefore given by:
$$\lambda (w_2|S|^4e^2)^{1/4}.$$
(80)
Close enough to the critical point, since $`|S|`$ becomes small, the skyrmion becomes large, and thus develops a topological character.
Now consider the excitations of the XY field. Taking $`\sigma _{i\mu }1`$ in Eq. (71), the naïve topological excitation consists of winding $`\theta `$ at infinity by an integer multiple of $`4\pi `$ – not $`2\pi `$, since the $`\mathrm{cos}(\mathrm{\Delta }_\mu \theta _i/2)`$ is not $`2\pi `$-periodic. These excitations are paired vortices in the 3-sublattice supersolid order parameter $`\psi `$. They are neutral, since there is no coupling to the dual gauge field. Like an ordinary neutral superfluid vortex, they cost a logarithmic energy, neglecting the XY anisotropy term $`w_3`$. When it is included, such a (double-strength) $`\psi `$ vortex becomes linearly confined, and converts at long distances to an intersection point of 12 (!) domain walls, the phase winding coalescing into these 12 walls radiating outward from the “vortex” core.
Allowing for a non-trivial texture in the $`Z_2`$ gauge field, however, a third kind of excitation is possible in the Mott phase. In particular, one may consider a “vison” or $`Z_2`$ vortex, around a point around which any line product of Ising gauge fields gives $`1`$. This requires the existence of a “cut”, a ray emanating outward from the vison along which Ising gauge fields crossing the ray are taken negative (nevertheless, the $`Z_2`$ flux is non-trivial only through one plaquette at the vison core). Such a cut effectively introduces anti-periodic boundary conditions for both $`z_\sigma `$ and $`e^{i\theta /2}`$ across the cut. That such a configuration is possible is of course evident from the definition of $`z_\sigma `$ and $`\theta `$ in Eq. (61), since the apparent discontinuities in $`z_\sigma `$ and $`e^{\pm i\theta /2}`$ do not affect the elementary $`z_{\pm \sigma }`$ fields. The anti-periodic boundary condition forces topological defects into both the XY and CP<sup>1</sup> fields. In the XY sector, it requires the existence of a $`\pm 2\pi `$ vortex in $`\psi `$. As for the $`\pm 4\pi `$ vortex above, this costs logarithmic energy neglecting $`w_3`$, and degenerates into a linearly confined “source” for (in this case 6) radial domain walls. In the CP<sup>1</sup> sector, it is slightly less intuitive. One might have expected the occurrence of some sort of “half-skyrmion” pseudospin texture. However, this is not the case. Anti-periodic boundary conditions require the discontinuity to persist all the way down from infinity to the vison core. There is thus no way to “relax” the winding singularity at infinity into a smooth pseudospin texture. Instead, the minimal energy configuration has $`𝐒`$ spatially constant everywhere, i.e. has the form $`z_\sigma =\xi _\sigma e^{i\mathrm{\Theta }}`$, where $`\mathrm{\Theta }`$ winds by $`\pm \pi `$ and $`\xi _\sigma `$ is constant everywhere save in some small core region of microscopic size. However, the continuum action Eq. (74) still obtains at infinity, so that finite energy configurations still satisfy Eq. (76) and $`A_\mu =𝒞_\mu `$ at infinity. Hence, these CP<sup>1</sup> “half-vortex” configurations carry fractional boson charges $`\pm 1/2`$. So by taking into account $`Z_2`$ gauge vortices, we find a third class of “elementary” excitations in the Mott insulator, “half bosons” with a texture in the 3-sublattice supersolid order parameter $`\psi `$. These are linearly confined beyond some length at which the 6-fold XY anisotropy becomes significant.
Of the three types of topological excitations discussed, it is interesting to note that only the charge $`\pm 1`$ skyrmion remains unconfined at the longest scales.
#### IV.3.3 SS3 phase
Let us now turn to the SS3 phase, which is described by Eq. (71) at large $`t_\theta `$. As discussed above, the ground state in this limit can be regarded as a $`Z_2`$ Higgs phase, with the $`z_\sigma `$ field uncondensed. Hence there are two types of topological defects: “paired” XY vortices in $`\psi `$, in which $`\theta `$ winds by a multiple of $`4\pi `$ and single XY vortices, in which $`\psi `$ winds by $`\pm 2\pi `$, accompanied by a vison. Both cost logarithmic energy at short scales, crossing over to linear confinement as do similar excitations in the Mott state. Finally, there are the CP<sup>1</sup> “particles” created by the $`z_\sigma `$ field, which can be considered to propagate coherently since the $`Z_2`$ gauge field is in a Higgs phase. The single XY vortices have a statistical interaction with the CP<sup>1</sup> quanta, but this does not lead to significant effects upon the CP<sup>1</sup> particles since the XY vortices are anyway linearly confined. The CP<sup>1</sup> quanta still carry unit dual gauge charge, and so should be regarded as the physical superfluid vortex excitations of the supersolid.
#### IV.3.4 SS2 Phase
The pseudospin vector order parameter $`𝐒=z^{}𝝉z^{}`$ is condensed in the SS2 phase. However, it is not a Higgs phase for the $`z_\sigma `$ fields, since, for instance, it is still superfluid, i.e. the dual gauge field remains gapless. Thus in the SS2 phase $`Z_2`$ quanta are strongly confined, and the elementary excitations must be $`Z_2`$ singlets. One class of excitations are skyrmions in $`𝐒`$. They do not carry any well-defined charge since the boson number conservation symmetry is anyway broken in the supersolid. The other quanta are physical vortices, which are bound states of $`z_\sigma `$ and $`e^{\pm i\theta /2}`$ particles, essentially the original $`z_{\pm \sigma }`$ vortices of the superfluid. However, because of the broken pseudospin symmetry in this phase, there is a preferred pseudospin polarization, and the two $`\sigma `$ components of the vortex spinor (choosing a quantization axis along $`𝐒`$) are no longer energetically equivalent. Therefore there is only a two-fold rather than four-fold low-energy vortex multiplet $`\xi _\pm z_{\pm 0}`$, taking $`\sigma =0`$ as the lower-energy spinor.
### IV.4 Supersolid-Mott quantum critical points
#### IV.4.1 SS3-Mott transition
The supersolid to Mott insulator QCPs are manifestly not of LGW type, if they occur at all as continuous transitions. Nevertheless, at least some aspects of these QCPs can be straightforwardly analyzed by application of Eq. (71). Consider first the SS3-Mott transition. It is useful to approach the transition first from the SS3 phase. Since it can be regarded as a $`Z_2`$ Higgs state, this is particularly simple. In particular, we can treat $`\sigma _{i\mu }1`$ as a constant at low energies. The $`e^{\pm i\theta _i/2}`$ fields can be regarded as condensed, and moreover at low energies there are no associated gapless Goldstone modes due to the 6-fold anisotropy term $`w_3`$ in Eq. (55). Hence at low energies, the only important fluctuations are those of the CP<sup>1</sup> spinor and the dual gauge field. The natural critical theory, kept lattice regularized for simplicity, and neglecting for the moment pseudospin anisotropy, is thus just
$`_{NCCP1}`$ $`=`$ $`t\sigma _{i\mu }e^{iA_{i\mu }}z_{i\sigma }^{}z_{i+\mu \sigma }+{\displaystyle \frac{1}{2e^2}}(ϵ_{\mu \nu \lambda }\mathrm{\Delta }_\nu A_{i\lambda })^2.`$
This is the Non-Compact CP<sup>1</sup> (NCCP<sup>1</sup>) theory studied numerically in Ref. Lesik, , and believed to represent a distinct universality class of critical phenomena. It also has a more intuitive interpretation: it describes the behavior of the $`2+1=3`$ dimensional $`O(3)`$ transition associated with $`𝐒`$ when “hedgehog” defects in $`𝐒`$ are completely suppressed in the partition function. Let us now see how this is understood approaching the QCP from the Mott side. In this phase, the pseudospin vector is ordered, but fluctuates more and more strongly as the SS3 phase is approached. One would expect skyrmion defects to become more and more prevalent as fluctuations in the pseudospin increase. As described above in Sec. IV.3.2, however, skyrmions in the Mott state carry physical boson charge, so boson number conservation requires that skyrmion-number ($`n_s`$ in Eq. (79)) must also be conserved. Happily, skyrmion number-changing events are exactly the hedgehog defects of the $`O(3)`$ model, so we see that charge conservation causes this transition to be of the NCCP<sup>1</sup> type.
Eq. (IV.4.1) and the subsequent discussion neglect pseudospin anisotropy. While it is quite likely such anisotropy terms are irrelevant the NCCP<sup>1</sup> fixed point, this requires further study. Using the hard-spin PSG transformations, Eqs. (163), the leading anisotropy terms can be shown to be
$`_{NCCP1}^{}`$ $`=`$ $`w_2{\displaystyle \underset{a=x,y,z}{}}(S^a)^4+w_2^{}\left(\mathrm{Im}\psi ^3\right)S^xS^yS^z.`$ (82)
The latter term is allowed by symmetry, but vanishes for $`w_3>0`$, in which case $`\psi ^3`$ is purely real. Note that the perturbations $`w_2,w_2^{}`$ in Eq. (82) are $`8^{\mathrm{th}}`$ and $`6^{\mathrm{th}}`$ order in the CP<sup>1</sup> fields, respectively, so it is quite plausible that both are irrelevant at the NCCP<sup>1</sup> point, though clearly this is most likely for $`w_3>0`$, when the $`w_2^{}`$ term is absent.
A further complication in the case $`w_3<0`$ is that the non-vanishing $`\psi `$ order parameter in this case breaks particle-hole symmetry $`C`$ (actually the supersolid with $`w_3>0`$ also breaks C, but preserves the combination $`CT_2R_{2\pi /3}I_{d_1}I_{d_2}`$, which is sufficient). Since a supersolid, like a superfluid, is compressible, this has the difficulty that it implies a non-vanishing deviation of the spatially averaged density from half-filling (working at fixed chemical potential chosen to maintain particle-hole symmetry – i.e. zero Zeeman field in the XXZ model). In the canonical ensemble, fixing the average density at $`f=1/2`$, it implies phase separation. Formally, this is described in the dual theory by the allowed coupling term
$$^{\prime \prime }=\lambda \left(\mathrm{Im}\psi ^3\right)(\mathrm{\Delta }_xA_y\mathrm{\Delta }_yA_x),$$
(83)
since the physical density is the dual magnetic flux. This indeed leads to a density deviation from half-filling away from the NCCP<sup>1</sup> critical point, since it leads to a minimum of $`_{NCCP1}+^{\prime \prime }`$ with non-zero $`\delta n(\mathrm{\Delta }_xA_y\mathrm{\Delta }_yA_x)/2\pi `$. As the Mott transition is approached, however, the system becomes increasingly less compressible, and the compressibility certainly vanishes when superfluidity does. It is not entirely clear to us how this is resolved – the complications are similar to (but more difficult due to the pseudospin structure) those occuring in the theory of the normal-superconducting thermal phase transition in a three-dimensional superconductor in a weak applied external field $`H`$.NelsonSeung It is possible that, at fixed chemical potential, the NCCP<sup>1</sup> critical fixed point is “weakly avoided” at long scales by this effect, most likely by introducing a narrow region of an “SS6” phase – a supersolid with the same symmetry as the Mott insulator but with ODLRO – between the Mott insulator and ferrimagnetic SS3 state. Working at fixed density $`f=1/2`$, one expects to pass through the NCCP<sup>1</sup> point, which coincides with the critical endpoint of the phase separation region.
#### IV.4.2 SS2 to Mott transition
As indicated in Sec. IV.3.4, the SS2 phase should be thought of as a state in which $`𝐒0`$, but vortices themselves are not condensed. It is not, however, a Higgs phase of Eq. (71). Moreover, the important elementary excitations of this phase are just vortices, which are bound states of the hard-spin fields. Therefore it is advantageous to return to the original soft-spin vortex formulation, and procede by just adding the term
$$_{SS2}^{}=\lambda 𝐒(𝐒_++𝐒_{}),$$
(84)
where of course the $`𝐒_\pm `$ fields on the right should be understood as composites of $`z_{\pm \alpha }`$. As discussed in Sec. IV.3.4, this splits the 4-fold vortex multiplet into two 2-fold multiplets. Taking $`𝐒`$ along $`(100)`$ (we will not discuss the $`(111)`$ case in any detail, but it is similar), we obtain the scalar low energy fields $`z_\pm `$, defined by $`z_{\pm \alpha }=z_\pm \eta _\alpha `$, with $`\tau ^x\eta =+\eta `$ (other orientations are solved by the obvious generalization). By considering the residual symmetry operations of this SS2 state (see Appendix D), we may thereby derive the continuum action for these two fields:
$`={\displaystyle \underset{\alpha =\pm }{}}\left[|(_\mu iA_\mu )z_\alpha |^2+s|z_\alpha |^2\right]+u(z_\alpha ^{}z_\alpha ^{})^2`$ (85)
$`v|z_+|^2|z_{}|^2+\lambda \mathrm{Im}(z_+^{}z_{}^{})^6+{\displaystyle \frac{1}{2e^2}}(ϵ_{\mu \nu \lambda }_\nu A_\lambda )^2.`$
Here we have kept the $`12^{\mathrm{th}}`$ order $`\lambda `$ term because it is the lowest order term which breaks the staggered $`U(1)`$ symmetry.
Remarkably, Eq. (85) is extremely similar to the NCCP<sup>1</sup> Lagrangian, differing mainly in that the SU(2) symmetry of that theory is here reduced by the $`v`$ term to U(1). For $`v>0`$, corresponding to the easy-plane case, it is the continuum theory for the “deconfined quantum critical point” of Refs. dcprefs, that describes the superfluid to VBS transition on the square lattice, with the modification that the “clock” anisotropy $`\lambda `$ is here 6-fold rather than 4-fold. This theory, neglecting the irrelevant $`\lambda `$ term, is self-dual at the critical point, and can alternatively be formulated as a theory of the bose condensation of charge $`\pm 1/2`$ fractional bosons. These are just the half-boson excitations discussed in Sec. IV.3.2 on the Mott state, which carry a direct U(1) gauge charge.
On examining the charge ordering pattern in Fig. 8, an interesting question arises. The symmetry of the $`(100)`$ SS2 state is already consistent with a very simple half-filled Mott insulator, consisting of stripes of charge on alternate lines of sites (along principle axes of the triangular lattice). So it is perfectly conceivable that in some models, one could have a transition from the SS2 supersolid to a Mott insulator with the same symmetry as the SS2 phase. One would expect this to be an XY transition, since only the superfluid $`U(1)`$ symmetry is broken across the transition. Why does our theory not predict this simpler scenario?
Firstly, we note that deconfined quantum criticality for the SS2 to Mott insulator transition studied above is perfectly consistent, since the Mott insulator in question is not the one with the same symmetry as the SS2 phase. The question remains why we do not see that possibility as well. Our interpretation is that, by starting with the dual field theory for the $`z_{\pm \sigma }`$ vortices, we have chosen a restricted set of vortex modes (this particular multiplet), which describes the natural instabilities of an isotropic triangular lattice superfluid. Though we have lowered the symmetry already in the SS2 phase, we have presumed the low energy excitations in this phase should be taken same vortex multiplet. If the spatial symmetry breaking present in the SS2 phase were taken strong, this might not be a good assumption, and states origination from other vortex multiplets could cross the $`z_\pm `$ states in energy, and lead to instabilities to different Mott states and also different critical behavior. We leave an exploration of this idea for future work.
#### IV.4.3 SS2 to spin liquid transition
On passing from the SS2 phase to the $`Z_2`$ spin liquid in Fig. 9, a vortex pair field $`z_\sigma z_\sigma ^{}`$ must condense. Because in the SS2 phase, the pseudospin rotational symmetry is already broken, we expect the pair field composed of two spinors aligned along the $`𝐒`$ axis to describe the condensate. This is a one-component field with dual U(1) gauge charge $`2`$. Hence we expect this transition to be described by a massless charge scalar coupled to a non-compact $`U(1)`$ gauge field. This is just dual to the XY model, so this can be viewed as an XY transition. In more physical terms, on passing from the $`Z_2`$ phase to the SS2 phase, the half-boson excitations of the spin liquid condense. It is clear from this description that the symmetry of the spin liquid is the same as that of the SS2 phase.
#### IV.4.4 Spin liquid to Mott transition
In the $`Z_2`$ spin liquid, both $`𝐒`$ and the vortex pair field are condensed. It has, however, the symmetry of the SS2 phase. It can be viewed as the Higgs phase of the $`\sigma _{i\mu }`$ gauge field. To pass to the Mott state, which does not have topological order, vison excitations must condense. The $`e^{\pm i\theta /2}`$ particles play the role of the vison. This can be seen, e.g. from the fact that these particles have a statistical interaction with the half-bosons in the $`Z_2`$ phase, which are $`\pi `$-flux tubes in $`A_{i\mu }`$. The non-standard feature is that the visons also carry space group quantum numbers – since the $`e^{\pm i\theta /2}`$ operator is a “square root” of the SS3 order parameter (the $`\theta `$ transformations are given in Appendix C). Actually this transition can be understood from Eq. (71) simply by “freezing” $`\sigma _{i\mu }1`$ and treating the $`z_\sigma `$ fields as condensed. It is thus clearlyan XY transition, with either $`6`$-fold or $`12`$-fold “clock” anisotropy (it is doubled since “half” the SS3 order parameter is condensing), for $`𝐒`$ along $`(100)`$ and $`(111)`$, respectively.
### IV.5 Anti-parallel pseudospins
Here we briefly sketch the results of an analogous study of the antiparallel pseudospin case. Using the parameterization in Eq. (III.3.3) to define the hard-spin degrees of freedom, and gauging the $`Z_2`$ redundancy, the appropriate hard spin model is
$`_{\mathrm{Z}_2}^{}`$ $`=`$ $`t_z\sigma _{i\mu }z_{i\sigma }^{}z_{i+\mu \sigma }t_\theta \sigma _{i\mu }\mathrm{cos}(\mathrm{\Delta }_\mu \theta _i/2A_{i\mu })`$ (86)
$`+{\displaystyle \frac{1}{2e^2}}(ϵ_{\mu \nu \lambda }\mathrm{\Delta }_\nu A_{i\lambda })^2.`$
Note that, in this case, the dual gauge field is coupled to the XY and not the CP<sup>1</sup> degree of freedom.
Study of various limits and interpolation determines the phase structure of Eq. (86), which is schematically shown in Fig. 10. In addition to the superfluid and Mott states occuring when both $`t_z,t_\theta `$ are small and large, respectively, there is a supersolid phase, SS12 (it has a minimal unit cell of $`12`$ sites), and a $`Z_2`$ spin liquid insulating phase. Unlike in the parallel pseudospin model, the SS12 supersolid has the same symmetry as the nearby Mott insulator. This is a direct consequence of the fact that the dual gauge charge is in this case carried by the $`U(1)`$ degree of freedom $`e^{i\theta }`$, and this does not carry any space group quantum numbers. For the same reason, the $`Z_2`$ spin liquid phase at large $`t_\theta `$ does not break any symmetries. Note that of the four phases in Fig. 10, only the $`Z_2`$ spin liquid is “exotic”, i.e. has an underlying topological order and unconventional excitations not captured by any local mean-field theory and order parameters.
The excitations are as follows. In the Mott state there are neutral skyrmions (since the $`z_\sigma `$ fields carry no dual gauge charge), confined charge $`\pm 1/2`$ excitations (visons bound to half-vortices in $`z_\sigma ,\theta `$), and charge $`\pm 1`$ excitations that can be viewed as $`\pm 4\pi `$ vortices in $`\theta `$. Because $`\theta `$ is coupled to $`A_{i\mu }`$, the latter cost finite energy, and are clearly adiabatically connected to vacancy/interstitials in the Mott state. The SS12 phase has the same skyrmion textures, but no well-defined charge excitations since it is a superfluid. Instead, it has a single physical vortex/antivortex excitation created by $`e^{\pm i\theta /2}`$. The “triviality” of the vortex multiplet is consistent with the broken symmetry of the supersolid, whose enlarged unit cell contains on average an integer number ($`6`$ in the simplest case) bosons. The $`Z_2`$ spin liquid has physical boson charge $`\pm 1/2`$ excitations ($`2\pi `$ vortices in $`\theta `$ accompanied by a “vison” in $`\sigma _{i\mu }`$) and physical “vison” excitations (created by $`e^{\pm i\theta /2}`$ and $`z_\sigma `$) which carry spatial quantum numbers.
The direct transition from superfluid to antiparallel Mott state is described by the continuum action of the previous section. The superfluid-SS12 and the SS12-Mott critical points are also conventional, since each is characterized by a change in a single order parameter. The superfluid-SS12 transition is described by an LGW theory for the $`𝐝`$-vector, while the SS12-Mott transition is simply an XY transition for the superfluid order parameter. The superfluid-$`Z_2`$ transition is also an XY transition, which can be understood as a condensation of charge $`\pm 1/2`$ “half-bosons” (in principle this changes universal amplitudes from the conventional superfluid to integer-filling Mott transition, which is also XY-likeSedgewick ). The $`Z_2`$ to Mott transition is modeled by the CP<sup>1</sup> action with no gauge field, which has the physical interpretation of modeling vison condensation. We have not attempted to consider the effects of various anisotropies on these transitions.
## V Discussion
In this paper we have presented a phenomenological dual vortex theory of the interplay between Mott localization and geometrical frustration for interacting bosons at half-filling on the triangular lattice. This approach reveals a variety of novel quantum phases and phase transitions which may occur if the superfluid and Mott insulating states occur in close proximity to one another in phase space. Of particular interest are the continuous superfluid-Mott insulator transition predicted by mean field theory, the two supersolid phases, and the occurrence of the recently-discovered NCCP<sup>1</sup> critical universitality class at the 3-sublattice supersolid to Mott insulator transition. In this discussion, we will provide a more direct physical picture of some of these phenomena, and address the prospects of observing them in simple microscopic boson or spin models.
A useful starting point for the discussion is the recent demonstration that a supersolid phase indeed occurs in the simplest spin-$`1/2`$ XXZ model,
$$H_{XXZ}=\underset{ij}{}J_{}(S_i^xS_j^x+S_i^yS_j^y)+J_zS_i^zS_j^z,$$
(87)
with ferromagnetic XY and antiferromagnetic Ising exchanges ($`J_{},J_z>0`$) (equivalently, hard-core bosons with nearest neighbor repulsion) on nearest-neighbor links of the triangular lattice.Melko05 ; Damle05 ; Troyer05 This model was shown to be in a 3-sublattice SS3-type phase for $`J_z5J_{}`$, and this phase persists up to and including $`J_z=\mathrm{}`$. A number of features of the numerical results on the supersolid at large $`J_z`$ are notable. First, although superfluidity survives, it is extremely weak, as characterized by the superfluid stiffness, which is approximately $`250`$ times smaller in the large $`J_z`$ supersolid than in the pure XY model ($`J_z=0`$). Second, it is exceedingly difficult to distinguish numerically on even relatively large lattices between the two different types of SS3 charge ordering patterns. Ref. Melko05, was unable to distinguish them numerically by direct measurement of boson density correlation functions, while Ref. Damle05, claimed to do so, but only for large lattices of $`18\times 18`$ sites with a very small signal. Furthermore, a deviation of the density from half-filling is expectedMelko05 in the “ferrimagnetic” SS3 phase identified in Ref. Damle05, (which corresponds in our theory to the one with $`w_3<0`$), but appears to be exceedingly minute if observable at all computationally. Apparently there is very little splitting energetically between the two SS3 states, even at $`J_z=\mathrm{}`$, a point at which there is no intrinsic small parameter in the microscopic Hamiltonian – the effective Hamiltonian is simply the XY exchange projected into the Hilbert space spanned by the manifold of classical Ising antiferromagnetic ground states on the triangular lattice.
The present theory offers a partial explanation for these puzzling observations. We interpret the weakness of superfluidity as evidence that the system is in close proximity to a Mott insulating state. If so, our dual vortex field theory, which is built around a superfluid to Mott insulating transition, should apply. The tiny energy splitting between the two SS3 states is then understandable: the term which dictates this splitting, $`w_3`$ is $`12^{\mathrm{th}}`$ order in the basic $`z_{\pm \sigma }`$ vortex fields, and clearly strongly irrelevant at the superfluid-Mott QCP. Furthermore, the smallness of any spontaneous density deviation from $`1/2`$-filling, which as indicated above is expected for the “ferrimagnetic” SS3 state, is also expected, since the density deviation must vanish as the Mott state, which has density of exactly $`1/2`$ and is incompressible, is approached. This deviation also vanishes at the transition from this state to the superfluid (as seen in LGW theoryMelko05 ), so it is likely small throughout the ferrimagnetic phase.
Given these arguments for proximity to the Mott state, it seems likely that only a small perturbation of the XXZ Hamiltonian may be required to push it into a Mott phase, and in so doing observe the very interesting NCCP<sup>1</sup> criticality at the SS3-Mott quantum critical point. Let us try to develop a more physical picture of this transition. We will begin by providing a cartoon understanding of the SS3 phases. In the case $`w_3<0`$, as discussed in Refs. Damle05, ; Troyer05, , the ferrimagnetic SS3 state can be understood crudely by first forming a “solid” of bosons with a density of $`1/3`$, with one boson occupying each of the sites of one of the three $`\sqrt{3}\times \sqrt{3}`$ triangular sublattices of the original lattice. This leaves, at $`f=1/2`$, a density of $`1/4`$ boson per the remaining sites, which form a honeycomb lattice, with the “solid” bosons in the centers of the honeycomb plaquettes. The SS3 state can be viewed as a superfluid of these remaining bosons (alternatively, one can make the same construction with holes replacing bosons, leading to a different but equivalent state – which illustrates the spontaneously broken particle-hole symmetry of the ferrimagnet). In the opposite case $`w_3>0`$, there is no spontaneous density polarization, and the “sublattice magnetizations” take the values $`S_i^z=(m,m,0)`$ on the three inequivalent sites. We will call this the “antiferromagnetic” state because of the exactly opposite Ising moments on two of the three sublattices. This state can be understood by a similar cartoon. In particular, take a honeycomb sublattice of the triangular lattice, and on this sublattice form a $`1/2`$-filled Mott insulator of alternating empty and occupied sites (an Ising Néel state in spin language). The sites at the centers of the honeycomb plaquettes form a $`\sqrt{3}\times \sqrt{3}`$ triangular sublattice, and we put the remaining bosons into a half-filled superfluid on this sublattice. These cartoon pictures would clearly tend to favor the ferrimagnetic state in the XXZ model, since bosons must hop (presumably virtually) between second neighbor sites to stabilize the triangular superfluid in the antiferromagnetic state. This immediately suggests that, to study the antiferromagnetic supersolid, one needs only to add a second neighbor XY exchange to the XXZ model,
$$H^{}=J_{}^{}\underset{ij}{}S_i^xS_j^x+S_i^yS_j^y.$$
(88)
With these pictures of the SS3 phases in hand, it is natural to view the transition to the Mott insulator as a “crystallization” of the superfluid sublattices of the supersolid. Indeed, we find that this provides a simple physical picture consistent with the symmetries of the appropriate Mott phases. Consider first the transition from the ferrimagnet. Here we have a $`1/4`$-filled honeycomb lattice of bosons, which is to undergo a superfluid to Mott insulator transition. At $`1/4`$-filling, these bosons are unlikely to form a simple “crystalline” Mott insulator, since they would be too widely separated to substantially interact. Instead, more likely Mott states are valence bond solids, in which the bosons resonate between two or more sites (within still-localized wavefunctions). The most natural candidate is a “columnar” valence bond solid (VBS) state, in which alternating columns of bonds are occupied by one boson. Indeed, the state predicted by the vortex mean field theory for $`w_3<0`$ and $`w_2<0`$ has exactly the symmetries of the columnar valence bond solid, see Fig. 11. One may also convince oneself of the validity of the columnar VBS picture by counting the number of distinct Mott states. Fix the location of the $`\sqrt{3}\times \sqrt{3}`$ superlattice and hence the honeycomb sublattice. One can then place the valence bonds along columns parallel to any of the three principle axes (of the honeycomb, which has principle axes halfway between those of the underlying triangular lattice), and for each such orientation, they may lie on even or odd columns. Hence one expects six states. This is precisely the number of distinct choices of vector pseudospin along the $`\pm \widehat{x},\pm \widehat{y},\pm \widehat{z}`$ axes.
Next consider the antiferromagnetic supersolid. We have an effective $`\sqrt{3}\times \sqrt{3}`$ triangular sublattice of bosons at half-filling, living in the plaquette centers of an “antiferromagnetic” bose solid on the honeycomb lattice. In the absence of this surrounding solid, the triangular sublattice would have, scaled up to its size, all the same symmetries as the original triangular lattice. The formation of a Mott insulator on this triangular sublattice would thus naïvely appear to be just as formidable a problem as the original one. The staggered solid, however, breaks $`I_{d_2}`$ and $`C`$, preserving only the combination, $`I_{d_2}C`$. This means, were we to repeat the PSG analysis for a dual vortex theory of the $`\sqrt{3}\times \sqrt{3}`$ sublattice bosons, there is no symmetry to prevent the Ising order parameter $`\mathrm{\Phi }=|z_+|^2|z_{}|^2`$ from appearing as a term in the action, breaking the $`\alpha =\pm `$ “flavor” degeneracy of the vortex multiplet. Clearly in such a theory, then only one of the two flavors will condense, and only those phases in which one pseudospin is non-zero will appear. Incidentally, one may identify therefore the $`z_\sigma `$ spinor in the hard-spin action with the lower-energy vortex flavor in this sublattice theory. The phases with only a single pseudospin condensed are “staggered” phases of the original bosons. The simplest of these is just a “columnar crystal”, in which bosons on alternating columns of triangular lattice sites (along some principle axis) are occupied and unoccupied. Superimposing this upon the surrounding antiferromagnetic honeycomb lattice, one remarkably obtains an ordering pattern again identical in symmetry to the original mean-field solid with $`w_3>0`$ and $`w_2<0`$! The counting of states is also the same as for the ferrimagnetic case above, since the “columns” have the same set of orientations and have only shifted from bonds to sites, once again in agreement with the configurations of $`𝐒`$.
One can easily extend these constructions to the cases with pseudospin along (111). For brevity, we relegate this to Appendix E.
The cartoons of the ferrimagnetic and antiferromagnetic supersolids/solids not only reproduce the symmetries of the phases, but also can be directly used to obtain the critical NCCP<sup>1</sup> theory of the SS3-Mott transition. In the ferrimagnetic case, the cartoon picture is a superfluid-Mott transition on a $`1/4`$-filled honeycomb lattice of bosons. In Appendix F, we sketch how a dual vortex analysis of that “cartoon” problem directly recovers the appropriate NCCP<sup>1</sup> theory for this case. In the antiferromagnetic case, a similar argument has already been sketched for the cartoon of an half-filled triangular $`\sqrt{3}\times \sqrt{3}`$ sublattice in the background of the honeycomb antiferromagnetic Mott state, which removes one of the two pseudospinors from the theory, likewise recovering the NCCP<sup>1</sup> model and the appropriate anisotropy.
It thus seems of considerable interest to pursue small perturbations of the XXZ model which might take it into a Mott phase, and between ferrimagnetic and antiferromagnetic states. To make the latter transition, from the ferrimagnetic to antiferromagnetic SS3 supersolid, we can make a simple recommendation: the inclusion of weak second-neighbor ferromagnetic exchange, as in Eq. (88), should favor the antiferromagnetic state by allowing bosons to better move on the triangular sublattice. The antiferromagnetic states seem slightly more favorable than the ferrimagnetic ones theoretically as candidates for observing the NCCP<sup>1</sup> criticality, because the latter preserve particle-hole symmetry in the supersolid, and have weaker pseudospin anisotropy terms. Understanding what terms microscopically would be needed to push the XXZ model toward the insulators is a nontrivial problem. However, it seems clear that introducing further-neighbor Ising exchange will lead to some Mott insulating states, when $`J_z`$ is large, since they spoil the degeneracy of the classical Ising manifold. Exploration of this problem numerically is tempting.
Analytically, it is quite interesting how naturally the dual vortex approach leads to the observed supersolid states. Coming from the Mott insulator with parallel pseudospins, the deconfined charge excitations are skyrmions. If one imagines lowering the charge gap in the insulator, it is most natural that these should condense, destroying the pseudospin order of the insulator and simultaneously initiating superfluidity. While other scenarios such as the direct Mott insulator to superfluid transition are possible theoretically (and described analytically in this paper), they entail a non-trivial screening so that confinement of other (fractionally charged) excitations of the Mott state becomes progressively weaker as the transition is approached. Thus the SS3 supersolid is probably the most intuitive “neighbor” of the Mott state, and with this understanding its occurence in the strong interaction limit of the XXZ model is no longer surprising. It is interesting to speculate as to whether it might be possible to formulate the boson-vortex duality used heavily in this paper directly in this strong-coupling region, i.e. in the classical Ising antiferromagnetic Hilbert space. Such a formulation could potentially provide a quantitative method for attacking problems of Mott transitions with strong frustration.
###### Acknowledgements.
We acknowledge useful discussions with O.I. Motrunich, A. Paramekanti and A. Vishwanath. Financial support was provided by the National Science Foundation grant DMR-9985255 and by the Packard Foundation.
## Appendix A Various PSG details
The PSG transformations in momentum space are given by:
$`T_1`$ $`:`$ $`\{\begin{array}{ccc}\psi _1(k_1,k_2)& & \psi _1(k_1,k_22\pi f)e^{ik_1}\hfill \\ \psi _2(k_1,k_2)& & \psi _2(k_1,k_22\pi f)e^{ik_1}\omega \hfill \end{array},`$ (91)
$`T_2`$ $`:`$ $`\{\begin{array}{ccc}\psi _1(k_1,k_2)& & \psi _1(k_1,k_2)e^{ik_2}\hfill \\ \psi _2(k_1,k_2)& & \psi _2(k_1,k_2)e^{ik_2}\hfill \end{array},`$ (94)
$`R_{2\pi /3}`$ $`:`$ $`\{\begin{array}{ccc}\psi _1(k_1,k_2)& & \frac{1}{q}_{m,n=0}^{q1}\omega ^{m(m+2n+\nu _q)/2}\psi _1(k_2+2\pi fm,k_1k_2+2\pi f(n1/2+\nu _q/2))\hfill \\ \psi _2(k_1,k_2)& & \frac{1}{q}e^{i(k_1+k_2)}_{m,n=0}^{q1}\omega ^{(m1)(m+2n1+\nu _q)/2}\psi _2(k_2+2\pi fm,k_1k_2+2\pi f(n1/2+\nu _q/2))\hfill \end{array},`$ (97)
$`I_{d_1}`$ $`:`$ $`\{\begin{array}{ccc}\psi _1(k_1,k_2)& & \frac{1}{q}_{m,n=0}^{q1}\psi _1^{}(k_2+2\pi fm,k_1+2\pi fn)\omega ^{mn}\hfill \\ \psi _2(k_1,k_2)& & \frac{1}{q}e^{i(k_1+k_2)}_{m,n=0}^{q1}\psi _2^{}(k_2+2\pi fm,k_1+2\pi fn)\omega ^{(m1)n}\hfill \end{array},`$ (100)
$`I_{d_2}`$ $`:`$ $`\{\begin{array}{ccc}\psi _1(k_1,k_2)& & \frac{1}{q}e^{ik_1}_{m,n=0}^{q1}\psi _2^{}(k_2+2\pi fm,k_1+2\pi fn)\omega ^{(m1)n}\hfill \\ \psi _2(k_1,k_2)& & \frac{1}{q}e^{ik_2}_{m,n=0}^{q1}\psi _1^{}(k_2+2\pi fm,k_1+2\pi fn)\omega ^{mn}\hfill \end{array},`$ (103)
where $`\nu _q=q2[q/2]`$.
Using these general forms for the PSG in momentum space, it is straightforward to obtain Eqs. (II,II) of the main text. The remaining PSG transformations for the multiplets are:
$`R_{2\pi /3}:\phi _{\mathrm{}}{\displaystyle \frac{1}{\sqrt{q}}}{\displaystyle \underset{\mathrm{}^{}=0}{\overset{q1}{}}}\omega ^{\mathrm{}(\mathrm{}+2\mathrm{}^{}+1)/2}\phi _{\mathrm{}^{}},`$
$`I_{d_1}:\phi _{\mathrm{}}{\displaystyle \frac{1}{\sqrt{q}}}{\displaystyle \underset{\mathrm{}^{}=0}{\overset{q1}{}}}\omega ^{\mathrm{}\mathrm{}^{}}\phi _{\mathrm{}^{}}^{},`$
$`I_{d_2}:\phi _{\mathrm{}}{\displaystyle \frac{1}{\sqrt{q}}}{\displaystyle \underset{\mathrm{}^{}=0}{\overset{q1}{}}}\omega ^{\mathrm{}\mathrm{}^{}}\phi _{\mathrm{}^{}}^{},`$ (104)
for $`q`$ odd, and
$`R_{2\pi /3}:\phi _{\alpha \sigma }e^{i\pi \alpha \sigma /q}{\displaystyle \frac{e^{i\eta _1(\alpha ,f)}}{\sqrt{q}}}{\displaystyle \underset{\sigma ^{}=0}{\overset{q1}{}}}\omega ^{\sigma (\sigma +2\sigma ^{}+1)/2}\phi _{\alpha \sigma ^{}},`$
$`I_{d_1}:\phi _{\alpha \sigma }{\displaystyle \frac{e^{i\eta _2(\alpha ,f)}}{\sqrt{q}}}{\displaystyle \underset{\sigma ^{}=0}{\overset{q1}{}}}\omega ^{\sigma \sigma ^{}}\phi _{\alpha \sigma ^{}}^{},`$
$`I_{d_2}:\phi _{\alpha \sigma }{\displaystyle \frac{e^{i\eta _3(\alpha ,f)}}{\sqrt{q}}}{\displaystyle \underset{\sigma ^{}=0}{\overset{q1}{}}}\omega ^{\sigma \sigma ^{}}\phi _{\alpha ,\sigma ^{}}^{},`$ (105)
when $`q`$ is an even integer. Appying $`R_{2\pi /3}`$ to $`\stackrel{~}{\varrho }_{mn}^\alpha `$, one finds that
$$\eta _1(\alpha )=\frac{\pi \alpha }{6q}.$$
(106)
Applying $`I_{d_2}`$, one finds that $`\eta _3(\alpha )=0`$. The remaining phase factor $`\eta _2(\alpha ,f)`$ is difficult to find analytically in the general case. In the case $`f=1/2`$ considered in detail in the text, it is, however, equal to $`\eta _1(\alpha )`$, so collecting these results,
$$\eta _1(\alpha )=\eta _2(\alpha )=\pi \alpha /12,\eta _3(\alpha )=0\text{for }f=1/2.$$
(107)
It is possible that $`\eta _1(\alpha )=\eta _2(\alpha )`$ holds at a general filling, not just at $`f=1/2`$, but it is not obvious.
For the specific case of $`f=1/2`$ ($`p=1,q=2`$), we can impose invariance under the particle-hole transformation,
$$C:\phi _{\pm \sigma }\phi _\sigma ^{}.$$
(108)
Imposing invariance under the rotation and reflection operations above upon the quartic action $`_1`$ taken in the form of Eq. (28) yields the set of conditions
$`\gamma _{mn}=\gamma _{m,n},`$
$`\gamma _{mn}=\gamma _{m,mn},`$
$`\gamma _{mn}=\gamma _{m2n,n}^{},`$
$`\gamma _{m^{}n^{}}={\displaystyle \frac{1}{q}}{\displaystyle \underset{m,n=0}{\overset{q1}{}}}\gamma _{mn}\omega ^{mn^{}nm^{}},`$
$`\gamma _{m^{}n^{}}={\displaystyle \frac{1}{q}}{\displaystyle \underset{m,n=0}{\overset{q1}{}}}\gamma _{mn}\omega ^{n(m^{}n2n^{})+m(n+n^{})},`$
for $`q`$ odd, and
$`\gamma _{mn}^{\alpha \beta }=\gamma _{m,n}^{\alpha \beta },`$
$`\gamma _{mn}^{\alpha \beta }=\gamma _{m,n}^{\beta \alpha },`$
$`\gamma _{mn}^{\alpha \beta }=\gamma _{m2n,n}^{\alpha \beta },`$
$`\gamma _{m^{}n^{}}^{\alpha \beta }={\displaystyle \frac{1}{q}}{\displaystyle \underset{m,n=0}{\overset{q1}{}}}\gamma _{mn}^{\alpha \beta }\omega ^{mn^{}nm^{}},`$
$`\gamma _{m^{}n^{}}^{\alpha \beta }={\displaystyle \frac{1}{q}}{\displaystyle \underset{m,n=0}{\overset{q1}{}}}\gamma _{mn}^{\alpha \beta }e^{i\pi n(\beta \alpha )/q}\omega ^{n(m^{}n2n^{})+m(n+n^{})},`$
for $`q`$ even.
## Appendix B Pseudospin transformations at $`f=1/2`$
For convenience, we give the transformation properties of the spinor $`z_{\alpha \sigma }`$ vortex fields. With the pseudospin index suppressed, we find
$`T_1`$ $`:`$ $`\{\begin{array}{ccc}z_+\hfill & & \hfill e^{i\pi /6}\tau ^xz_+\\ z_{}\hfill & & \hfill e^{i\pi /6}\tau ^xz_{}\end{array},`$ (113)
$`T_2`$ $`:`$ $`\{\begin{array}{ccc}z_+\hfill & & \hfill e^{i\pi /6}\tau ^zz_+\\ z_{}\hfill & & \hfill e^{i\pi /6}\tau ^zz_{}\end{array},`$ (116)
$`R_{2\pi /3}`$ $`:`$ $`\{\begin{array}{ccc}z_+\hfill & & \hfill e^{i\pi /3}u_rz_+\\ z_{}\hfill & & \hfill e^{i\pi /3}u_rz_{}\end{array},`$ (119)
$`I_{d_1}`$ $`:`$ $`\{\begin{array}{ccc}z_+\hfill & & \hfill e^{i5\pi /12}u_1^{}z_+^{}\\ z_{}\hfill & & \hfill e^{i5\pi /12}u_1^{}z_{}^{}\end{array},`$ (122)
$`I_{d_2}`$ $`:`$ $`\{\begin{array}{ccc}z_+\hfill & & \hfill iu_2^{}z_{}^{}\\ z_{}\hfill & & \hfill iu_2^{}z_+^{}\end{array},`$ (125)
$`C`$ $`:`$ $`\{\begin{array}{ccc}z_+\hfill & & \hfill ϵz_{}^{}\\ z_{}\hfill & & \hfill ϵz_+^{}\end{array}.`$ (128)
Here
$`u_r`$ $`=`$ $`e^{i\frac{\pi }{3}\widehat{𝒏}_r𝝉}={\displaystyle \frac{1}{2}}\left(\begin{array}{cc}1+i& 1+i\\ 1+i& 1i\end{array}\right),`$ (131)
$`u_1`$ $`=`$ $`e^{i\frac{\pi }{2}\widehat{𝒏}_1𝝉}={\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{cc}i& i\\ i& i\end{array}\right),`$ (134)
$`u_2`$ $`=`$ $`e^{i\frac{\pi }{2}\widehat{𝒏}_2𝝉}={\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{cc}i& i\\ i& i\end{array}\right),`$ (137)
with
$`\widehat{𝒏}_r`$ $`=`$ $`(1,1,1)/\sqrt{3},`$
$`\widehat{𝒏}_1`$ $`=`$ $`(1,0,1)/\sqrt{2},`$
$`\widehat{𝒏}_2`$ $`=`$ $`(1,0,1)/\sqrt{2}.`$ (138)
Next we tabulate the transformation properties of the pseudospin vector order parameters:
$`T_1`$ $`:`$ $`\{\begin{array}{ccc}S_\alpha ^x\hfill & & \hfill S_\alpha ^x\\ S_\alpha ^y\hfill & & \hfill S_\alpha ^y\\ S_\alpha ^z\hfill & & \hfill S_\alpha ^z\end{array},`$ (142)
$`T_2`$ $`:`$ $`\{\begin{array}{ccc}S_\alpha ^x\hfill & & \hfill S_\alpha ^x\\ S_\alpha ^y\hfill & & \hfill S_\alpha ^y\\ S_\alpha ^z\hfill & & \hfill S_\alpha ^z\end{array},`$ (146)
$`R_{2\pi /3}`$ $`:`$ $`\{\begin{array}{ccc}S_\alpha ^x\hfill & & \hfill S_\alpha ^y\\ S_\alpha ^y\hfill & & \hfill S_\alpha ^z\\ S_\alpha ^z\hfill & & \hfill S_\alpha ^x\end{array},`$ (150)
$`I_{d_1}`$ $`:`$ $`\{\begin{array}{ccc}S_\alpha ^x\hfill & & \hfill S_\alpha ^z\\ S_\alpha ^y\hfill & & \hfill S_\alpha ^y\\ S_\alpha ^z\hfill & & \hfill S_\alpha ^x\end{array},`$ (154)
$`I_{d_2}`$ $`:`$ $`\{\begin{array}{ccc}S_\alpha ^x\hfill & & \hfill S_\alpha ^z\\ S_\alpha ^y\hfill & & \hfill S_\alpha ^y\\ S_\alpha ^z\hfill & & \hfill S_\alpha ^z\end{array},`$ (158)
$`C`$ $`:`$ $`𝐒_\alpha 𝐒_\alpha .`$ (159)
## Appendix C Hard spin tranformations
Here we give the transformation rules for the hard-spin theory with parallel pseudospins. Defining
$$z_\pm =ze^{\pm i\theta /2},$$
(160)
we obtain
$`T_1`$ $`:`$ $`\{\begin{array}{ccc}z\hfill & & \hfill i\tau ^xz\\ \theta \hfill & & \hfill \theta 4\pi /3\end{array},`$ (163)
$`T_2`$ $`:`$ $`\{\begin{array}{ccc}z\hfill & & \hfill i\tau ^zz\\ \theta \hfill & & \hfill \theta 4\pi /3\end{array},`$ (166)
$`R_{2\pi /3}`$ $`:`$ $`\{\begin{array}{ccc}z\hfill & & \hfill u_rz\\ \theta \hfill & & \hfill \theta 2\pi /3\end{array},`$ (169)
$`I_{d_1}`$ $`:`$ $`\{\begin{array}{ccc}z\hfill & & \hfill u_1^{}z^{}\\ \theta \hfill & & \hfill 5\pi /6\theta \end{array},`$ (172)
$`I_{d_2}`$ $`:`$ $`\{\begin{array}{ccc}z\hfill & & \hfill u_2^{}z^{}\\ \theta \hfill & & \hfill \theta \end{array},`$ (175)
$`C`$ $`:`$ $`\{\begin{array}{ccc}z\hfill & & \hfill ϵz^{}\\ \theta \hfill & & \hfill \theta +\pi .\end{array}`$ (178)
The $`SU(2)`$ matrices $`u`$ are defined in Eqs. (131). We have used the freedom to redefine any of these operations by a global $`U(1)`$ gauge transformation, which in this hard-spin limit corresponds to a phase rotation of $`z`$.
## Appendix D Residual symmetries of the SS2 supersolid
The residual symmetry operations of this SS2 state (which are not broken by the pseudospin vector order) are generated by $`T_1`$, $`R_{2\pi /3}I_{d_1}`$, $`CI_{d_2}I_{d_1}`$, and $`T_2I_{d_2}I_{d_1}`$. Their action on the $`z_\pm `$ vortex fields can be obtained from the definition $`z_{\pm \sigma }=z_\pm \eta _\sigma `$ and Eqs. (113):
$`T_1`$ $`:`$ $`z_\pm \pm e^{i\pi /6}z_\pm ,`$
$`R_{2\pi /3}I_{d_1}`$ $`:`$ $`\{\begin{array}{ccc}z_+& & z_+^{}\\ z_{}& & iz_{}^{}\end{array}`$ (182)
$`CI_{d_2}I_{d_1}`$ $`:`$ $`z_\pm e^{i\pi /12}z_\pm ^{},`$
$`T_2I_{d_2}I_{d_1}`$ $`:`$ $`z_\pm e^{\pm i\pi /12}z_{}.`$ (183)
These lead directly to Eq. (85).
## Appendix E Toy model wavefunctions with (111) pseudospin
Here we provide toy model wavefunctions for the Mott states with $`w_2>0`$ and parallel pseudospins, so that $`𝐒_+=𝐒_{}`$ lies along say the $`(111)`$ axis. First consider $`w_3>0`$, in which we wish to construct a state appropriate to a $`1/4`$-filled honeycomb lattice. In this case, a wavefunction with the correct symmetry is shown in Fig.13. In this figure the bosons denoted by squares are on the honeycomb sublattice of the triangular lattice, and become superfluid in the SS3 phase.
Now consider $`w_3<0`$, in which case we must consider the half-filled triangular lattice embedded in the antiferromagnetic honeycomb Mott insulator. The density plot for this state is shown in Fig.14. Focusing on the triangular lattice sites, we note that the amplitude on these sites has the same $`2\times 2`$ periodicity as the $`w_3>0`$ state above, with the sites on this sublattice having one amplitude, and the remaining sites another. However, it is inconsistent with half-filling to simply put one boson on the sublattice sites, or on the other three sites. Moreover, the state is rather constrained by the existence of a number of centers of three-fold rotations and reflections. The simplest Mott wavefunction we constructed that satisfies all symmetry properties is actually not quite a product state of the form of Eq. (39). It is, however, clearly an insulating state. To construct it, assume the $`2\times 2`$ sublattice sites are empty. Then the remaining sites form a kagome lattice. The kagome lattice is composed of corner-sharing triangles, or two orientations (pointing “left” and “right” in the figure). We act with creation operators to place one boson on each triangle in a uniform superposition.
$$\mathrm{\Psi }=\underset{\mathrm{}}{}(b_\mathrm{}1^{}+b_\mathrm{}2^{}+b_\mathrm{}3^{})|0$$
(184)
Since the triangles overlap, there will actually be amplitude to find more than one boson per triangle. On average, the number of bosons per unit cell of the kagome lattice is then $`2`$, since there is one up and one down triangle per unit cell. There are four triangular lattice sites for each such unit cell (3 kagome and one central empty site), so this is a state at $`1/2`$-filling, which manifestly has all the symmetries of the Mott state in Fig.14. Moreover, it is an insulating state as required, since one can readily show there are only very short-range correlations, e.g. in the boson Green’s function.
## Appendix F Honeycomb lattice at $`1/4`$-filling
The vortex PSG’s on the honeycomb lattice were worked out in Appendix B of the first paper in Ref. Balents04, . They are sufficient to develop a dual vortex theory of the superfluid-Mott transition on this lattice. For the case of $`f=1/4`$, there are two vortex flavors, $`\phi _0,\phi _1`$, which can be combined into a spinor $`\phi _\sigma `$. We consider the restrictions upon the vortex Lagrangian by translations and rotations. Suppressing the indices and using Pauli matrices $`𝝉_{\sigma \sigma ^{}}`$ to span the spinor space, we have:
$`T_1`$ $`:`$ $`\phi \tau ^x\phi ,`$
$`T_2`$ $`:`$ $`\phi \tau ^z\phi ,`$
$`T_d`$ $`:`$ $`\phi \tau ^y\phi ,`$
$`R_{2\pi /3}^{\mathrm{dual}}`$ $`:`$ $`\phi e^{i\pi /3}u_r\phi ,`$ (185)
where $`u_r`$ is given in Eq. (131). Requiring invariance under these PSG operations, one may readily write down a continuum Lagrangian in terms of $`\phi `$:
$`_{\mathrm{honey}}`$ $`=`$ $`{\displaystyle \underset{\sigma =0,1}{}}\left[|(_\mu iA_\mu )\phi _\sigma |^2+s|\phi _\sigma |^2\right]+u(\phi _\sigma ^{}\phi _\sigma ^{})^2`$ (186)
$`+`$ $`{\displaystyle \frac{1}{2e^2}}(ϵ_{\mu \nu \lambda }_\nu A_\lambda )^2+_a,`$
where $`SU(2)`$ symmetry of $`\phi _\sigma `$ is first broken by the term
$$_a=\lambda S^xS^yS^z,$$
(187)
with $`𝐒=\phi _\sigma ^{}𝝉_{\sigma \sigma ^{}}\phi _\sigma `$. Eq. (186) indeed has the global $`SU(2)`$ and $`U(1)`$ gauge symmetry of the NCCP<sup>1</sup> model, and $`_a`$ recovers the leading anisotropy term in Eq. (82) for the ferrimagnetic case.
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# Spontaneous Parity Violation
## I Quantum Loop Operators
The partition function, $`𝒵`$, for a general continuous group $`G`$ we define in terms of quantum loop operators by,
$$𝒵_N(t)=dg\mathrm{exp}\left[_0^tA_g(𝒏_s)zV(𝒏_s)ds\right],$$
(1)
with $`t,z`$,
$`A_g(𝒏)`$ $``$ $`{\displaystyle \underset{(i,j)}{\overset{NT}{}}}{\displaystyle \underset{\sigma G}{}}\lambda _{ij\sigma }(𝒏){\displaystyle \frac{𝒏\mathrm{𝟏}_{i\sigma }\mathrm{𝟏}_{j\sigma }|g}{𝒏|g}},`$ (2)
$`V(𝒏)`$ $``$ $`{\displaystyle \underset{(i,j)}{\overset{NT}{}}}{\displaystyle \underset{\sigma G}{}}\lambda _{ij\sigma }^{}(𝒏){\displaystyle \frac{𝒏\mathrm{𝟏}_{i\sigma }\mathrm{𝟏}_{\sigma j}|𝒏}{𝒏|𝒏}},`$ (3)
where $`\lambda _{ij\sigma }(𝒏),\lambda _{ij\sigma }^{}(𝒏)`$, $`\{\sigma \}`$ form some finite subset of the elements of $`G`$, and $`g`$ is a general element of $`G`$. We treat the complex-valued continuous symmetry generalisation of the problem considered in presilla-gauge ; presilla-old . The fact that the operator definition is diagonalisable implicitly defines a local gauge transformation on the Hilbert space,
$`{\displaystyle \underset{(i,j)}{\overset{NT}{}}}`$ $`{\displaystyle \underset{\sigma =G}{}}\lambda _{ij\sigma }(𝒏){\displaystyle \frac{𝒏^{}\mathrm{𝟏}_{i\sigma }\mathrm{𝟏}_{j\sigma }|𝒏}{𝒏|g}}𝒏^{}|g`$ (4)
$`=`$ $`𝒏^{}|g𝒏^{}|A|𝒏𝒏|g^1=𝒏^{}|A_g|𝒏.`$
We should compare this directly with the single plaquette model given in terms of the Vandermonde determinant discussed in vandermondea ; vandermondeb , which yields a similar form of time-ordered product. The connection to Vafa-Witten formalism is made explicit by noticing that equivalently the integration of the exponent in can be performed over the spatial Lattice volume via a change of basis, $`𝒏_s^{}`$. The partition function is then given in the usual Vafa-Witten form,
$`𝒵_T(z)`$ $`=`$ $`{\displaystyle dg\mathrm{exp}\left[_0^zA_g(𝒏_s^{})tV(𝒏_s^{})ds^{}\right]}`$ (5)
$`=`$ $`{\displaystyle dg\mathrm{exp}\left[_0^zA_g^{}(𝒏_s^{})itV^{}(𝒏_s^{})ds^{}\right]}`$
$`A_g^{}(𝒏_s^{})`$ $`=`$ $`\mathrm{Im}[A_g(𝒏_s^{})],`$ (6)
$`V^{}(𝒏_s^{})`$ $`=`$ $`V(𝒏_s^{})\mathrm{Re}[A_g(𝒏_s^{})],`$ (7)
where $`V(𝒏_s^{})`$ since it is hermitean in the previous basis, $`𝒏_s`$.
### I.1 Thermodynamic Limit Monotonicity
It is surprisingly simple to prove the existence of the thermodynamic limit for partition function in (1). The proof consists of two parts. The first step is to identify the saddle point solution of the partition function. The second step is then to determine the limiting value of this asymptotic solution when the neighbourhood of the saddle point is extended to cover the entire phase space. Apart from isolated singularities the partition function is then necessarily holomorphic. This is the essence of Lee and Yang’s discussion in lee+yang ; poissona .
For the quantum case we face now an additional subtlety. This has not been previously addressed in the literature. Partition functions for quantum systems are defined over an extended phase space - a trace is taken over an additional dimension. Two examples are given above in (1) and (5). For a quantum system to be truly holomorphic, and the thermodynamic limit exist, it is therefore not merely sufficient that the thermodynamic limit be proven to exist over either $`𝒏_s`$ or $`𝒏_s^{}`$. Rather, the thermodynamic limit must exist over $`𝒏`$. Expressing this another way, existing partition function studies invoking this holomorphic property over $`𝒏_s^{}`$ (holomorphic in $`z`$) have a thermodynamic limit that is strictly speaking only defined up to infrared divergences (ie. they are not necessarily simultaneously holomorhpic in $`t`$). For Lee and Yang with classical systems, which have no trace relation, this complication does not arise. The partition function in this case is given directly over $`𝒏`$, without the time-ordering of the integrand.
We now prove the existence of the thermodynamic limit for the new quantum loop operator formulation for the partition function in (1). The saddle point, which maximises the partition function in some region of the extended phase space $`M`$, is defined through,
$$R_M\underset{M}{sup}R(𝒎),$$
(8)
where $`𝒎`$ is parameterises some contour in $`NT`$, and $`R`$ is the integrand of the partition function. The value of the integrand in the neighbourhood of the saddle point, $`M_\delta `$, is defined through,
$$R_\delta (𝒎)R(𝒎)R_M0.$$
(9)
Not every point contained in $`M`$ is necessarily in the neighbourhood of the saddle point. We consider how the neighbourhood evolves as the contour is increased in length : the limit where $`𝒎𝒏`$.
$$\underset{𝒎𝒏}{lim}R_\delta (𝒎)=\underset{𝒎𝒏}{lim}R(𝒎)R_M0.$$
(10)
The relation between the consecutive basis elements (along $`𝒎`$) is given partly through the local gauge symmetry transformation in (4). However, we can have a situation for some general group $`G`$ in which we have multiple degenerate solutions of (10). In this case a subspace must be projected out to time-order the integrand and resolve this degeneracy. With this resolution is then necessarily monotonic and decreasing, since the potential nonlocal ambiguities are removed. Next, taking the limit $`lim_{MNT}`$ we have,
$`\underset{MNT}{lim}`$ $`[\underset{𝒎𝒏}{lim}R_\delta (𝒎)]`$ (11)
$`=`$ $`\underset{MNT}{lim}[\underset{𝒎𝒏}{lim}R(𝒎)R_M]0.`$
This limit is also monotonic and decreasing with above choices, and therefore the thermodynamic limit of the quantum loop operator partition function in is well-defined. The holomorphic properties we are therefore interested in are those of the space $`NT`$, modulo, the time-ordering necessary to resolve the nonlocal ambiguities of the gauge group. It should be clear now from a numerical perspective that it is not necessary to perform either integral in the integrands exactly : in order to be simultaneously holomorphic in $`t`$ and $`z`$ it is sufficient for the integrand to cover $`NT`$.
### I.2 Saddle Point Determination
We are interested to find an explicit expression for the evolution of the partition function, $`𝒵_N`$, in the limit $`N\mathrm{},N>0`$. There are exactly $`N`$ singularities in the partition function in as a consequence of the above monotonicity arguments. By construction, the loop operator definition of $`𝒵_N`$ is meromorphic so we can equally well formulate the discussion in terms of the poles of $`𝒵_N`$, $`\{V_k\}`$, or the zeroes, $`\{\lambda _k\}`$. The poles turn out to be easier to treat than the zeroes, but for clarity the relation between the poles and zeroes is given through,
$`A(z)`$ $``$ $`{\displaystyle \underset{k=0}{\overset{N}{}}}{\displaystyle \frac{z+\lambda _k}{z+V_k}},`$ (12)
$`{\displaystyle \frac{1}{2\pi i}}{\displaystyle _𝒞}{\displaystyle \frac{A^{}(z)}{A(z)}}𝑑z`$ $`=`$ $`{\displaystyle \underset{k}{}}n(𝒞,\lambda _k){\displaystyle \underset{k}{}}n^{}(𝒞,V_k),`$ (13)
where $`n`$ and $`n^{}`$ are the respective homotopies of the zeroes and poles, and $`𝒞`$ is some contour homologous to $`0`$ in $`NT`$. This property will become important shortly when we consider finite multiplicities for the singularities, or equivalently, global symmetries in $`NT`$.
The schema we now follow is the complex-valued generalisation of the problem considered in presilla-saddle ; bessel . Namely, we define the Laplace transform of the partition function, $`\stackrel{~}{𝒵}_N`$, as a function of the poles, $`\{V_k\}`$. Then we formulate the inverse Laplace transform that gives us back the partition function expressed as a function of $`t`$, $`N`$ and $`\{V_k\}`$. Finally, we evaluate this integral equation for $`𝒵_N`$ via standard asymptotic series approaches in order to identify the dominant behaviour of the partition function in the above limit. In particular, this approach allows us to identify the saddle point solution of the partition function as distinct from the asymptotic convergence properties of the zeroes towards the $`\mathrm{Re}z`$ axis.
The Laplace transform of the partition function polynomial in at finite $`N`$ is given in terms of the poles of the partition function by,
$$\stackrel{~}{𝒵}_N(z)=_0^{\mathrm{}}𝑑te^{zt}𝒵_N(t)=ϵ^N\underset{k=0}{\overset{N}{}}\frac{1}{z+V_k},$$
(14)
where $`z,ϵ,V_k`$. The inverse Laplace transform of is given by,
$$𝒵_N(t)=\frac{1}{2\pi i}_𝒞𝑑ze^{zt}\stackrel{~}{𝒵}_N(z)=\frac{1}{2\pi iϵ}_𝒞𝑑z\mathrm{exp}\left[N\phi (z)\right],$$
(15)
$`\phi (z)={\displaystyle \frac{zt}{N}}{\displaystyle \underset{k=0}{\overset{N}{}}}{\displaystyle \frac{1}{N}}\mathrm{log}\left({\displaystyle \frac{z+V_k}{ϵ}}\right).`$ (16)
The contour $`𝒞`$ in, is defined as any suitable path running between $`\pm \mathrm{i}\mathrm{}`$, for which any singular behaviour is appropriately treated. We want to now evaluate the asymptotic behaviour of the partition function separately for the two cases : when a parity symmetry is, and is not, realised on the vacuum.
In the first case we have no special symmetry relations defined between the poles, and $`\phi (z)`$ is therefore of the generic form above in. Solving explicitly for $`\phi ^{}(z_0)=0`$ we have,
$$\underset{k=0}{\overset{N}{}}\frac{1}{z_0+V_k}=t.$$
(17)
Note, the location of the saddle point, $`z_0`$, is a function of both $`\{V_k\}`$ and $`t`$. Potentially the saddle point lies to the left of some subset of the poles, $`\{V_m\},mN`$. Therefore the integral in is given by,
$$𝒵_N(t)=\frac{1}{2\pi iϵ}_{𝒞_{z_0}}\mathrm{exp}\left[N\phi (z)\right]𝑑z+\underset{l=1}{\overset{m}{}}\frac{e^{V_lt}}{ϵ}\underset{kl}{\overset{N}{}}\frac{1}{V_l+V_k}$$
(18)
where $`𝒞_{z_0}`$ is the contour that passes through $`z_0`$. For this first case, when one assumes no parity symmetry is realised with the gauge fields, the asymptotic behaviour of the partition function is therefore given by,
$`\underset{N\mathrm{},N>0}{lim}𝒵_N(t)`$ $`=`$ $`{\displaystyle \frac{1}{2\pi iϵ}}\mathrm{exp}\left[N\phi (z_0){\displaystyle \frac{i}{2}}\theta (z_0)\right]\left(\sqrt{{\displaystyle \frac{2\pi }{N|\phi ^{\prime \prime }(z_0)|}}}+O(N^{3/2})\right),`$ (19)
$`=`$ $`\mathrm{exp}\left[z_0t{\displaystyle \underset{k=0}{\overset{N}{}}}\mathrm{log}\left({\displaystyle \frac{z_0+V_k}{ϵ}}\right){\displaystyle \frac{i}{2}}\theta (z_0)\right]\left(\sqrt{{\displaystyle \frac{1}{2\pi ϵ^2N|\phi ^{\prime \prime }(z_0)|}}}+O(N^{3/2})\right),`$ (20)
where $`\theta (z_0)=\mathrm{Arg}\phi ^{\prime \prime }(z_0)`$.
For the second case it is straightforward to define a parity symmetry action on $`Z_N(t)`$ simply by reversing the orientation of the elementary plaquettes in $`NT`$. This then yields the symmetry of the zeroes $`\{\lambda _k,\lambda _k^{}\}`$. From the Argument Principle in (13) the poles are then given through,
$$\stackrel{~}{𝒵}_N(z)=ϵ^N\underset{k=0}{\overset{N}{}}\frac{1}{\sqrt{(z+V_k)(z+V_k^{})}},$$
(21)
$$\phi (z)=\frac{zt}{N}\underset{k=0}{\overset{N}{}}\frac{1}{N}\mathrm{log}\left(\frac{\sqrt{(z+V_k)(z+V_k^{})}}{ϵ}\right).$$
(22)
The singularities at $`z=\mathrm{Re}V_k\pm \mathrm{Im}V_k`$ therefore correspond to branch points, with a branching exponent of $`1/2`$. Again, we want to define the contour such that it passes through $`z_0`$, but this is again potentially obstructed by singularities lying to the right of $`z_0`$. We therefore choose the branch cut to correspond to this deformation over the saddle point. If $`\mathrm{Im}V_k<0`$ the deformation will start at the branch point and descend to $`i\mathrm{}`$. If the $`\mathrm{Im}V_k>0`$ the deformation will start at the branch point and ascend to $`i\mathrm{}`$. If $`\mathrm{Im}V_k=0`$ the singularity is a pole, which we will treat as a residue contribution as before. This choice of branch cut does not introduce an arbitrariness into the calculation. Although it is not now possible to analytically continue across the cut, $`\phi (z)`$ is specified unambiguously along the original contour, and the analytic continuation between contours is also specified unambiguously. Crucially we keep the same branch cut for the pairs $`\{V_l,V_l^{}\}`$. The cuts to the left and right of the contour deformation $`𝒞`$ are defined through,
$$(z+V_l)_+^{1/2}=e^{\pi i}(z+V_l)_{}^{1/2}.$$
(23)
The analogue of the residue contribution for the branch points is therefore,
$`𝒵_N^l(t)`$ $`=`$ $`{\displaystyle _𝒞}e^{N\phi (z,kl)}(z+V_l^{})^{1/2}(z+V_l)_{}^{1/2}𝑑z`$
$``$ $`{\displaystyle _𝒞}e^{N\phi (z,kl)}(z+V_l^{})^{1/2}(z+V_l)_+^{1/2}𝑑z`$
$`=`$ $`\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}2}{\displaystyle _𝒞}e^{N\phi (z,kl)}(z+V_l^{})^{1/2}(z+V_l)_{}^{1/2}𝑑z`$
$$2e^{N\phi (z=V_l,kl)}\underset{n=0}{\overset{\mathrm{}}{}}\frac{\mathrm{\Gamma }(n+1/2)b^n(0)}{n!N^{n+1/2}}$$
(26)
where,
$$b(0)=(2i\mathrm{Im}V_l.\phi ^{}(V_l,kl))^{1/2},$$
(27)
$$\phi ^{}(z=V_l,kl)=\frac{t}{N}\underset{k=0,kl}{\overset{N}{}}\frac{1}{N}\frac{\mathrm{Re}(V_l+V_k)}{\sqrt{(V_l+V_k)(V_l+V_k^{})}}.$$
(28)
and the $`b^n(0)`$ form the analogue of the Laurent series for the branch points. The solutions are not oscillatory in $`N`$, as was found with previous discussions. In our new formulation we can see as well that the behaviour of the branch points can wholly dominate the asymptotic form of the partition function in the case in which a parity symmetry is realised, solely dependent on the relative magnitude of the prefactor, $`\mathrm{exp}[\phi (V_k,t)]`$.
We are able to reach this new point of unambiguous saddle point solutions for two reasons. Firstly, the new problem is defined over a higher dimension than previous discussions, so $`z`$ and $`t`$ are separated. This means that the singularities can be distinguished from asymptotic scaling. Secondly, because the new construction is meromorhpic we can identify all the singularities in $`𝒵_N`$, which are of a finite number.
The above arguments also are readily generalised to cases with different branching exponents - corresponding to different global symmetries. Specifically, if the branching exponent is less than -1 then the above asymptotic identifications for the branches can be performed via integration by parts.
## II spontaneous symmetry breaking
We have constructed in (1) a partition function that is formally equivalent to that of the Vafa-Witten discussion. It consists of a local bosonic path integral, on which is defined a symmetry breaking source term, $`A_g^{}(𝒏_s^{})`$. To relate our saddle point identification further to the Vafa-Witten discussion on spontaneous parity violation we now need to evaluate the vacuum expectation of the source term. Since the path integral is time-ordered this is given by,
$`𝒪`$ $`=`$ $`\underset{N\mathrm{}}{lim}\left({\displaystyle \frac{1}{N}}{\displaystyle \frac{\mathrm{log}\left(𝒵_N(t)\right)}{t}}\right)`$ (29)
$`=`$ $`\underset{N\mathrm{}}{lim}\left({\displaystyle \frac{1}{N}}{\displaystyle \frac{𝑑zze^{zt}\stackrel{~}{𝒵}_N(z)}{𝑑ze^{zt}\stackrel{~}{𝒵}_N(z)}}\right).`$ (30)
We therefore have two possible solutions, from (22) and (26), for the case when parity is realised by the system,
$$𝒪=z_0\mathrm{or}𝒪=V_l.$$
(31)
In particular we can see for this second case in (26), associated with the branch point, that the free energy density minima does not correspond with the partition function maxima at $`z_0`$. What we can see through our new formalism is that at finite $`N`$ there is a finite subtraction to made to properly define the contour for the partition function, and this itself can give rise to a parity breaking contribution. From (22) at most one of these expectations can be zero simultaneously, and so the potential for spontaneous symmetry breaking is well-defined.
Our new analysis does not mean that it is particularly likely that parity be spontaneously broken in QCD, since in our formulation it is identified as a quantum effect driven by the system dynamics through $`\{V_k\}`$. Plausibly, these symmetry breaking effects are only observed in instances where the statistical properties of the the Lattice ensemble become very much non quasi-classical, ie. the orders of the contributions in (22) and (26) become comparable hotb ; hotc .
As a second result we should notice from that it is possible to analytically continue the partition function throughout the entire $`z`$ plane. By definition the partition function is holomorphic in this region. We can deduce in principle it is therefore possible to resolve the complex action problem via numerical simulation. If a lattice determination correctly determines the polynomial expansion in up to quantum fluctuations at any single point in $`z`$, ie. that generates a covering of the Hilbert space in (2) and (3), the system can be smoothly mapped to all point in the $`z`$ plane. Since the result smoothly maps to the continuum the Lattice evaluation is also sufficiently defined for the dynamics to generate a nonzero vacuum expectation value of the source, $`𝒪`$, and to observe spontaneous symmetry breaking in the continuum.
Conversely, we can understand the origin of the disconnected contributions in existing Lattice evaluations. Algorithms can be based around the principle that a global symmetry is exactly held, and equal Monte Carlo weights assigned for each branch. However, because only the positive definite contributions are kept, corresponding to real probabilistic weights, the opposite cut is in effect taken for each branch point multiplet in (21). Genuine quantum effects can never therefore be observed as the cancellation between sectors is exact configuration by configuration, by construction and furthermore the system is nonrenormalizable StrongCP ; 4th .
## III Summary
We have introduced a quantum loop operator formalism for the partition function of vector-like gauge field theories, and identified that this is formally equivalent to the local bosonic operator expansion on finite system size introduced in the Vafa-Witten discussion. The Lattice formulation is given over an extended phase space $`NT`$, and the source parameter is itself discretized in the new approach. We have then considered the continuum limit of this space and identified that the zeroes of the loop operator polynomial are formally equivalent to partition function zeroes. Although this result is distinct from classical systems requiring a prior noncommuting limit to be taken over $`T`$ to ensure the formulation is meromorphic. The saddle point solution for the asymptotic behaviour of this complex-valued partition function was then determined. Spontaneous symmetry breaking is well-defined for case in which a parity symmetry is realised by the system. In the latter case the free energy minima can be distinct from the partition function maxima, and the generalisation of the result to general global symmetries was identified.
By construction, the formulation can be analytically continued across the entire phase plane in the zeroes expansion parameter. Therefore, any nonperturbative operator satisfying the quantum loop operator definition is an exact solution to the complex action problem.
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# DESY 05-096SFB/CPP-05-23MS-TP-05-15June 2005 Lattice Spacing Dependence of the First Order Phase Transition for Dynamical Twisted Mass Fermions
## 1 Introduction
Understanding the phase structure of lattice QCD is an important pre-requisite before starting large scale simulations. Indeed, our collaboration found that when working at lattice spacings of about $`0.15\mathrm{fm}`$ there can be strong first order phase transitions at small quark masses, at least when a combination of Wilson plaquette action and Wilson fermions is used . The phenomenon appears also when a small twisted mass term is switched on. This has serious consequences, since in such a scenario the pion mass $`m_{\mathrm{PS}}`$ cannot be made arbitrarily small but assumes a minimal value, $`m_{\mathrm{PS}}^{\mathrm{min}}`$, which may be about $`500\mathrm{MeV}`$ and hence it becomes impossible to work close to the physical value of the pion mass.
The presence of the first order phase transition for pure Wilson fermions is in accordance with predictions from chiral perturbation theory , which have been extended later to the case of adding a twisted mass . Let us, for completeness, also mention that for values of the lattice spacing much coarser than $`a=0.15\mathrm{fm}`$ the first order phase transition turns into a second order one from the normal QCD phase to the so-called Aoki phase . The generic phase structure of lattice QCD according to our present understanding is discussed and illustrated in refs. .
In refs. we have studied only one value of the inverse gauge coupling $`\beta =6/g_0^2`$ in order to demonstrate the existence of the first order phase transition, leaving the question of the $`\beta `$ dependence open. Since lattice chiral perturbation theory predicts a weakening of the first order phase transition towards the continuum limit, it is interesting to check this prediction and, in particular, to investigate quantitatively how fast the transition weakens when the continuum limit is approached. The answer to the latter question will naturally depend on the choice of the actions that are used for the gauge and the fermion fields.
In this paper we will present results using Wilson twisted mass fermions and the Wilson plaquette gauge action for three values of $`\beta `$. At each of these $`\beta `$ values we have performed simulations at a number of quark masses on both sides of the first order phase transition. This allows to study the $`\beta `$ dependence of the phase transition itself and, in addition, the lattice spacing dependence of physical observables computed separately in the two phases. We have performed such a scaling test for the pion mass, the pion decay constant and the ratio of the pion to the vector meson mass. For a scaling test of Wilson twisted mass fermions and other recent results in the quenched approximation see refs. .
## 2 Wilson twisted mass fermions
In this paper we will work with Wilson twisted mass fermions that can be arranged to be $`O(a)`$ improved without employing specific improvement terms . The Wilson tmQCD action in the twisted basis can be written as
$$S[U,\chi ,\overline{\chi }]=a^4\underset{x}{}\overline{\chi }(x)(D_W+m_0+i\mu \gamma _5\tau _3)\chi (x),$$
(1)
where the Wilson-Dirac operator $`D_\mathrm{W}`$ is given by
$$D_\mathrm{W}=\underset{\mu =0}{\overset{3}{}}\frac{1}{2}[\gamma _\mu (_\mu ^{}+_\mu )a_\mu ^{}_\mu ]$$
(2)
and $`_\mu `$ and $`_\mu ^{}`$ denote the usual forward and backward derivatives and the Wilson parameter $`r`$ was set to $`1`$.
The situation of full twist and hence automatic $`O(a)`$ improvement arises when $`m_0`$ in eq. (1) is tuned towards a critical bare quark mass $`m_{\mathrm{crit}}`$. We use for our simulations the hopping representation of the Wilson-Dirac operator with $`\kappa =(2am_0+8)^1`$.
We extract the pseudo scalar mass $`m_{\mathrm{PS}}`$ and the vector meson mass $`m_\mathrm{V}`$ from the usual correlation functions:
$$\begin{array}{cc}\hfill C_{PP}(x_0)& =a^3\underset{𝐱}{}P^+(x)P^{}(0),\hfill \\ \hfill C_{VV}(x_0)& =\frac{a^3}{3}\underset{k=1}{\overset{3}{}}\underset{𝐱}{}V_k^+(x)V_k^{}(0),\hfill \end{array}$$
(3)
where we consider the local bilinears $`P^\pm =\overline{\chi }\gamma _5\frac{\tau ^\pm }{2}\chi `$ and $`V_\mu ^\pm =\overline{\chi }\gamma _\mu \frac{\tau ^\pm }{2}\chi `$. Here we used $`\tau ^\pm =(\tau _1\pm i\tau _2)`$ with $`\tau _{1,2}`$ the first two Pauli matrices. Similarly one can define the correlation function $`C_{AP}`$ with the local bilinear $`A_\mu ^\pm =\overline{\chi }\gamma _\mu \gamma _5\frac{\tau ^\pm }{2}\chi `$.
The bare pseudo scalar decay constant $`f_\chi ^{\mathrm{PS}}`$ in the twisted basis can be obtained from (cf. )
$$f_\chi ^{\mathrm{PS}}=m_{\mathrm{PS}}^1r_{AP}0|P^+(0)|\pi ,$$
(4)
where the ratio
$$r_{AP}=\frac{0|A_0^+(0)|\pi }{0|P^+(0)|\pi }$$
(5)
can be extracted from the asymptotic behavior of
$$\frac{C_{AP}(x_0)}{C_{PP}(x_0)}=r_{AP}\mathrm{tanh}[m_{\mathrm{PS}}(T/2x_0)].$$
(6)
The bare PCAC quark mass $`m_\chi ^{\mathrm{PCAC}}`$ in the twisted basis can then be computed from the ratio
$$m_\chi ^{\mathrm{PCAC}}=\frac{f_\chi ^{\mathrm{PS}}}{20|P^+(0)|\pi }m_{\mathrm{PS}}^2.$$
(7)
The sign of $`m_\chi ^{\mathrm{PCAC}}`$ and $`f_\chi ^{\mathrm{PS}}`$ is determined by the sign of $`r_{AP}`$ and therefore, the corresponding values can be negative. One has to keep in mind that $`m_\chi ^{\mathrm{PCAC}}`$ and $`f_\chi ^{\mathrm{PS}}`$, since measured in the twisted basis, do not correspond to the physical quark mass and the physical pseudo scalar decay constant, respectively. While the quark mass is given by a combination of the (renormalized) values of $`m_\chi ^{\mathrm{PCAC}}`$ and $`\mu `$, the pseudo scalar decay constant can be computed by the help of $`f_\chi ^{\mathrm{PS}}`$ and the twist angle, as long as $`f_\chi ^{\mathrm{PS}}0`$ and the value of the twist angle is different from $`\pi /2`$.
Note that the purpose of the present paper is *not* to work at full twist nor to extract physical quantities, but rather to study the lattice spacing dependence of the first order phase transition. For the same reason, we also do not address the question of the choice of the critical quark mass in order to stay at full twist here, see refs. for recent quenched simulations addressing this point.
## 3 The phase transition as a function of the lattice spacing
In order to study the lattice spacing dependence of the phase transition we have chosen three values of $`\beta `$: $`\beta =5.1`$, $`\beta =5.2`$ and $`\beta =5.3`$. We scaled the volumes and the values of $`\mu `$ such that the physical volume is larger than $`2\mathrm{fm}`$, roughly constant and that $`r_0\mu 0.03`$, where $`r_0`$ is the Sommer scale fixed to be $`r_0=0.5\mathrm{fm}`$. Note that the value of $`r_0/a`$ depends on the value of the quark mass and therefore we had to choose a reference value for $`r_0/a`$ as will be explained below. The parameters are summarized in table 1.
In practice it turned out that a very direct way of detecting the presence of a first order phase transition in lattice QCD is to monitor the behavior of the plaquette expectation value $`P`$, e.g. as a function of $`\kappa `$ for fixed twisted mass parameter $`\mu `$. In such a situation, starting at identical parameter values from “hot” (random) or “cold” (ordered) configurations, $`P`$ can assume different, co-existing values. In fig. 1 we show $`P`$ as a function of $`1/(2\kappa )`$ for the three values of $`\beta `$. The picture is typical for the behavior of a first order phase transition with meta-stable branches, one with a low value of $`P`$ and one with a high value of $`P`$. We will denote in the following these branches as high (“H”) and low (“L”) plaquette phases, respectively.
The $`\beta `$-dependence shows that the gap in the plaquette expectation value $`\mathrm{\Delta }P`$ decreases substantially when moving from $`\beta =5.1`$ ($`a0.20\mathrm{fm}`$) to $`\beta =5.3`$ ($`a0.12\mathrm{fm}`$), which is presumably due to the mixing with the chiral condensate as discussed in . One possible definition for the quantity $`\mathrm{\Delta }P`$ is the difference between low and high phase plaquette expectation value at the smallest value of $`\kappa `$ where a meta-stability occurs.
Let us remark that the first order phase transition exists also in the continuum limit at zero quark mass where the scalar condensate has a jump as a consequence of spontaneous chiral symmetry breaking. This means, of course, that in the continuum limit the phase transition occurs only for $`\mu =0`$.
We give our simulation parameters, the statistics of the Monte Carlo runs and the results for $`am_{\mathrm{PS}}`$, $`af_\chi ^{\mathrm{PS}}`$, $`am_\chi ^{\mathrm{PCAC}}`$ and $`r_0/a`$ in tables 4, 5 and 6.
The meta-stability phenomenon observed in $`P`$ can also be seen in fermionic quantities. As an example, we show in fig. 2 the values of the PCAC quark mass as obtained in the branches with high and low plaquette expectation values of fig. 1 for the three values of $`\beta `$. Again we observe that with increasing $`\beta `$ the gap between positive (low plaquette phase) and negative (high plaquette phase) quark masses shrinks. Also, the meta-stability region in $`1/(2\kappa )`$ gets much narrower with increasing $`\beta `$.
The effects of the first order phase transition can also be seen in the pion mass and the value of the force parameter $`r_0`$. We plot in fig. 3 an example of the pion mass as a function of the PCAC quark mass at $`\beta =5.3`$. The most intriguing observation here is that due to the presence of the first order phase transition, the pion mass, say for positive quark masses, does not go to zero but rather reaches a minimal value, and jumps then to the phase with negative quark mass. This is, of course, just another manifestation of the jump in the PCAC quark mass in fig. 2.
In fig. 4 we also show the values of $`r_0/a`$ in the low and high plaquette phases at $`\beta =5.3`$. Note that the values of $`r_0/a`$ are quite different when determined in the low and the high plaquette phases, which is a generic feature also for other values of $`\beta `$ and even for different gauge actions, see ref. .
An interesting question is, at which value of the lattice spacing $`a`$ the minimal pion mass $`m_{\mathrm{PS}}^{\mathrm{min}}`$ assumes a value of, say, $`300\mathrm{MeV}`$ where contact to chiral perturbation could be established.
The pion mass assumes two different values for a fixed quark mass, once this quark mass lies inside the meta-stability region. These two values for the pion masses correspond to the two phases that for a certain interval of quark masses co-exist. The precise determination of the meta-stability region is, of course, very difficult. We can, however, give an interval in $`\kappa `$, $`[\kappa _1,\kappa _2]`$, that can be read from tables 4, 5 and 6 for the three different $`\beta `$ values, where meta-stabilities occur in our simulation. In the following, we will mainly concentrate on the low plaquette phase since this is the natural choice for studying lattice QCD. Being interested only in the low plaquette phase we determine then a lower bound for the minimal pion mass as computed at the lower end of this interval, i.e. $`\kappa _1`$, in the low plaquette phase. We give in table 2 the values of the minimal pion masses in the low plaquette phase in physical units. In addition, we provide the value for the gap in the plaquette expectation value $`\mathrm{\Delta }P`$.
In principle, it would be interesting to extrapolate the minimal pion mass and the gap in $`P`$ as a function of the lattice spacing. However, our present data do not allow for a reliable and safe extrapolation. First of all, the determination of the minimal pion mass has a large ambiguity in itself since we do not know exactly for which value of the quark mass the meta-stability will disappear. A substantially larger statistics would be necessary to answer this question and to check whether tunneling from one phase to the other occurs. Second, the only three values of $`\beta `$ we have used give a too short lever arm to perform a trustworthy extrapolation. And, third, the values of $`r_0/a`$ are very different in the two phases, as can be seen in fig. 4, which makes it particularly difficult to follow the gap in $`P`$ as a function of $`a/r_0`$.
Nevertheless, an estimate on a more qualitative level yields a value of the lattice spacing of $`a0.07\mathrm{fm}0.1\mathrm{fm}`$ where simulations with pion masses of about $`300\mathrm{MeV}`$ can be performed without being affected by the first order phase transition.
## 4 Lattice spacing dependence of <br>physical observables
Although the present simulations are not at full twist, the fact that we have results at three values of $`\beta `$ with roughly constant $`r_0\mu `$ allows us to check for the size of lattice artifacts. In order to perform such an investigation it is advantageous to express physical quantities in dimensionless variables. To this end, let us first define a reference pion mass through $`(r_0m_{\mathrm{PS}})^2=1.5`$. We have chosen this particular value in order to be able to interpolate for the values of $`\beta =5.1`$ and $`\beta =5.3`$, and to perform only a short extrapolation for $`\beta =5.2`$ to this point.
At the aforementioned reference pion mass, a corresponding reference value of $`r_0/a`$ and a reference quark mass can be determined, the latter leading to a variable $`\sigma `$,
$$\sigma =\frac{m_\chi ^{\mathrm{PCAC}}}{.m_\chi ^{\mathrm{PCAC}}|_{\mathrm{ref}}}.$$
(8)
Similarly, we can define ratios for a quantity $`O`$,
$$R_O=\frac{O}{.O|_{\mathrm{ref}}}$$
(9)
where $`O|_{\mathrm{ref}}`$ is the quantity as determined at the reference pion mass. The values for several quantities at the reference point can be found in table 3.
In order to determine the reference values for $`m_\chi ^{\mathrm{PCAC}}`$, $`f_\chi ^{\mathrm{PS}}`$ and $`r_0`$, in a first step we interpolated $`m_\chi ^{\mathrm{PCAC}}`$ linearly as a function of $`(r_0m_{\mathrm{PS}})^2`$ to the point where $`(r_0m_{\mathrm{PS}})^2=1.5`$ and extracted the reference value for $`m_\chi ^{\mathrm{PCAC}}`$. Then we determined the reference values for $`f_\chi ^{\mathrm{PS}}`$ and $`r_0`$ by quadratically interpolating the data as a function of $`m_\chi ^{\mathrm{PCAC}}`$ to the reference value of $`m_\chi ^{\mathrm{PCAC}}`$. We repeated the latter step with a linear interpolation finding agreement within the errors. The fits to the data have been performed with the ROOT and MINUIT packages from CERN (cf. ), taking the errors on both axis into account. We remark that for the quantity $`r_0m_{\mathrm{PS}}`$ we have neglected the correlation of the data between $`r_0/a`$ and $`am_{\mathrm{PS}}`$.
For a given observable $`O`$, $`R_O`$ is a universal function of $`\sigma `$ for fixed value of $`\mu `$ in physical units that allows for a direct comparison of results obtained at different values of $`\beta `$ and, in principle, even for different actions. Deviations of results at different $`\beta `$ values provide then a direct measure of scaling violations. In fig. 5 we show $`R_{m_{\mathrm{PS}}^2}`$ as a function of $`\sigma `$. Note that for the scaling analysis we take the data in the low plaquette phase only since this corresponds to the standard lattice QCD situation. We also remark that some of the points taken in this analysis might be meta-stable. Nevertheless, we assume here that these data can serve for checking scaling violations. Besides the data from the present work, we added also results from simulations at $`\beta =5.6`$ , which were obtained, however, at vanishing twisted mass parameter $`\mu =0`$.
A rather amazing consequence of fig. 5 is that, despite the fact that we are using coarse lattices, we cannot detect any scaling violation, at least within the (large) statistical errors of our data. Even more, the results of our present simulations at small values of $`\beta `$ agree with results from simulations with pure Wilson fermions at $`\beta =5.6`$ setting $`\mu =0`$. The same observation is made for $`R_{f_\chi ^{\mathrm{PS}}}`$, see fig. 6 and the ratio $`m_{\mathrm{PS}}/m_\mathrm{V}`$, see fig. 7. These results indicate that the lattice artifacts and the effect of a non-vanishing twisted mass parameter $`\mu `$ are surprisingly small. We remark here that in the case of the ratios like $`R_{m_{\mathrm{PS}}^2}`$ and $`R_{f_\chi ^{\mathrm{PS}}}`$ one could have cancellation of mass independent cutoff effects. One has also to have in mind that, due to the presence of the first order phase transition, the simulated pion masses are still larger than $`500\mathrm{MeV}`$. Whether our findings also hold when one is approaching the chiral limit is certainly an interesting but open question. However, in a set-up with Wilson twisted mass fermions and Wilson plaquette gauge action this question cannot be answered at these values of the lattice spacing.
## 5 Conclusions
In this paper we have investigated dynamical Wilson twisted mass fermions employing the Wilson plaquette gauge action. We have performed simulations at three values of $`\beta =5.1,5.2,5.3`$, corresponding to values of the lattice spacing of $`a0.20,0.16,0.14\mathrm{fm}`$, respectively. The non-zero values of the twisted mass parameter $`\mu `$ were chosen such that $`r_0\mu 0.03`$ for all of the three $`\beta `$ values. At these rather coarse lattice spacings we find clear signals of first order phase transitions that manifest themselves in a meta-stable behavior of the plaquette expectation value and fermionic quantities, such as the PCAC quark mass and the pion mass.
We clearly observe that the gaps in quantities sensitive to the phase transition, such as the plaquette expectation value and the PCAC quark mass decrease substantially when $`\beta `$ is increased. Unfortunately, with our present set of simulations, we are not able to quantitatively locate the value of the lattice spacing, where the effects of the first order phase transition becomes negligible and where a minimal pion mass of, say, $`300\mathrm{MeV}`$ can be reached. As an estimate of such a value of the lattice spacing we give a range of $`a0.07\mathrm{fm}0.1\mathrm{fm}`$. Of course, this would mean that a continuum extrapolation of physical results obtained on lattices with linear extent of at least $`L=2\mathrm{fm}`$ would be very demanding, since the starting point for such simulations would already require large lattices. It is therefore very important to find alternative actions such that the value of the lattice spacing can be lowered without running into problems with the first order phase transition. One candidate for such an action, the DBW2 gauge action, is discussed in ref. where it has indeed been found that modifying the gauge action alone can substantially reduce the strength of the first order phase transition. We are presently investigating another possibility, the tree-level Symanzik improved gauge action .
Despite the problems arising from the presence of a first order phase transition, we performed a scaling analysis for the pion mass, the pion decay constant and the ratio $`m_{\mathrm{PS}}/m_\mathrm{V}`$ for the data obtained at the three values of $`\beta `$ where we performed simulations. To this end, we only analyzed data from the low plaquette phase, since this is the natural choice for QCD simulations.
By defining a reference pion mass at $`(r_0m_{\mathrm{PS}})^2=1.5`$, we computed the ratio of $`m_{\mathrm{PS}}`$ and $`f_\chi ^{\mathrm{PS}}`$ to the corresponding reference values as a function of the PCAC quark mass, again measured with respect to the corresponding reference quark mass. We find that for these ratios the scaling violations are remarkably small and cannot be detected with the present precision of our data. Even more, when adding data from simulations of Wilson fermions with $`\mu =0`$ at $`\beta =5.6`$, then these data fall on the same scaling curve as our results on much coarser lattices and with twisted mass parameter switched on. This indicates that not only the lattice artifacts but also the effect of switching on a twisted mass of the order of $`r_0\mu 0.03`$ are small, at least for the rather large pion masses simulated here. This finding is surprising since it suggests that continuum values of physical quantities can be already estimated from simulations at not too small lattice spacings. Of course, our scaling results suffer from the fact that they are obtained using data that might be meta-stable as a consequence of the presence of the first order phase transition. Hence, a scaling test with an action that does not lead to significant effects of the first order phase transition is mandatory to check the results presented in this paper.
## 6 Acknowledgments
We thank R. Frezzotti, G. Münster, G. C. Rossi and S. Sharpe for many useful discussions. The computer centers at NIC/DESY Zeuthen, NIC at Forschungszentrum Jülich and HLRN provided the necessary technical help and computer resources. We are indebted to R. Hoffmann and J. Rolf for leaving us a MatLab program to check our fits. This work was supported by the DFG Sonderforschungsbereich/Transregio SFB/TR9-03.
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# LATTICE QCD AT FINITE TEMPERATURE
## 1 Introduction
It is expected that strongly interacting matter shows qualitatively new behavior at temperatures and/or densities which are comparable or larger than the typical hadronic scale. It has been argued that under such extreme conditions deconfinement of quarks and gluons should take place, i.e. thermodynamics of strongly interacting matter could be understood in terms of this elementary degrees of freedom and this new form of matter was called Quark Gluon Plasma $`^\mathrm{?}`$. On the lattice the existence of the deconfinement transition at finite temperature was first shown in the strong coupling limit of QCD $`^\mathrm{?}`$, followed by numerical Monte-Carlo studies of lattice SU(2) gauge theory which confirmed it $`^\mathrm{?}`$. Since these pioneering studies QCD at finite temperature became quite a large subfield of lattice QCD (for recent reviews on the subject see Refs. $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$). One of the obvious reasons for this is that phase transitions can be studied only non-perturbatively. But even at high temperatures the physics is non-perturbative beyond the length scales $`1/(g^2T)`$ ($`g^2(T)`$ being the gauge coupling) $`^\mathrm{?}`$. Therefore lattice QCD remains the only tool for theoretical understanding of the properties of strongly interacting matter under extreme condition which is important for the physics of the early universe as well as heavy ion collisions.
## 2 Finite temperature transition in full QCD
One of the basic questions we are interested in is what is the nature of the transition to the new state of mater and what is the temperature where it happens <sup>a</sup><sup>a</sup>aI will talk here about the QCD finite temperature transition irrespective whether it is a true phase transition or a crossover and $`T_c`$ will always refer to the corresponding temperature.. In the case of QCD without dynamical quarks, i.e. SU(3) gauge theory these questions have been answered. It is well established that the phase transition is 1st order $`^\mathrm{?}`$. Using standard and improved actions the corresponding transition temperature was estimated to be $`T_c/\sqrt{\sigma }=0.632(2)`$ $`^\mathrm{?}`$ ($`\sigma `$ is the string tension). The situation for QCD with dynamical quarks is much more difficult. Not only because the inclusion of dynamical quarks increases the computational costs by at least two orders of magnitude but also because the transition is very sensitive to the quark masses. Conventional lattice fermion formulations break global symmetries of continuum QCD (e.g. staggered fermion violate the flavor symmetry) which also introduces additional complications. Current lattice calculations suggest that transition in QCD for physical quark masses is not a true phase transition but a crossover $`^{\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?}}`$. Recent lattice results for the transition temperature $`T_c`$ from Wilson fermions $`^{\mathrm{?},\mathrm{?}}`$, improved $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$ and unimproved staggered fermions $`^\mathrm{?}`$ with 2 and 2+1 flavors of dynamical quarks are summarized in Fig. 1. The errors shown in Fig. 1 are only statistical with the exception of the data point from the MILC collaboration, where the large error partly comes from the continuum extrapolation and also includes systematic error in scale setting. Since the “critical” energy density $`ϵ_c=ϵ(T_c)`$ (i.e. the energy density at the transition) scales as $`T_c^4`$ the error in $`T_c`$ is the dominant source of error in $`ϵ_c`$ $`^\mathrm{?}`$.
## 3 Heavy quarks at finite temperature
In this section I am going to summarize some recent progress made in understanding the interaction of heavy quarks at finite temperature. Apart from being an interesting problem from a theoretical perspective understanding the interaction of heavy quarks at finite temperature also is very important for phenomenology. It has been suggested that quarkonium suppression due to color screening at high temperatures can serve as signature of Quark Gluon Plasma formation in heavy ion collisions $`^\mathrm{?}`$. For static quarks one can calculate the free energy difference for the system with static quark anti-quark pair and the system without static charges. This quantity is often referred to as finite temperature potential, though it should be emphasized that it is a free energy and thus contains an entropy contribution $`^\mathrm{?}`$. In Fig. 2 I show the free energy of static $`Q\overline{Q}`$ in the singlet state calculated in three flavor QCD $`^\mathrm{?}`$. As one can see from the Figure the free energy goes to a constant at distances $`r>0.9`$ fm at low temperatures. This happens because once enough energy is accumulated the string can break due to creation of a light quark-antiquark pair. As the temperature increases the distance where the free energy levels off becomes temperature dependent and decreases with increasing temperature. This reflects the onset of chromo-electric screening. Similar results have been obtained in two flavor QCD $`^{\mathrm{?},\mathrm{?}}`$.
It should be noticed that at short distances ($`r<0.4`$ fm) the free energy of static $`Q\overline{Q}`$ is temperature independent. As expected at short distance medium effects are not important. This is also confirmed by studies of the coupling constant at finite temperature $`^\mathrm{?}`$ which I also show in Fig. 2. The running of the coupling constant at finite temperature is controlled by the distance between the static quarks and its value is never larger than at zero temperature $`^\mathrm{?}`$. This disfavors the picture of strongly coupled plasma where $`\alpha _s`$ runs to large value above the transition temperature $`^\mathrm{?}`$.
Though the study of the free energy of a static quark anti-quark pair gives some useful insight into the problem of quarkonium binding at high temperatures ( for a recent review on this subject see Ref. $`^\mathrm{?}`$ ), it is not sufficient for detailed understanding quarkonium properties in this regime. To gain quantitative information on this problem quarkonium correlators and spectral functions should be studied at finite temperature. Such studies became possible only recently and still are restricted to the quenched approximation $`^{\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?}}`$. The results of these studies for charmonia are summarized in Fig. 3. The $`1S`$ states ($`J/\psi ,\eta _c`$) seem to survive to temperatures as high as $`1.6T_c`$ (maybe even higher, cf. the figure) while the 1P states ($`\chi _{c0},\chi _{c1}`$) are dissolved at $`1.1T_c`$ $`^\mathrm{?}`$. The survival of the $`1S`$ state is also confirmed by Umeda et al $`^\mathrm{?}`$. The temperature dependence of the charmonia correlators also suggests that the properties of $`1S`$ charmonia are not affected significantly above $`T_c`$, at least at zero spatial momentum $`^{\mathrm{?},\mathrm{?}}`$. As for bottomonia only preliminary results exists showing that $`\mathrm{{\rm Y}}`$ can exist in the plasma up to much higher temperatures $`^\mathrm{?}`$ but surprisingly enough the $`\chi _b`$ state is dissociated at temperatures smaller than $`1.5T_c`$ $`^\mathrm{?}`$.
## Acknowledgments
This work has been authored under the contract DE-AC02-98CH10886 with the U.S. Department of energy. I would like to thank Frithjof Karsch for careful reading of the manuscript.
## References
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# Unusual transport properties of ferromagnetic Heusler alloy Co2TiSn
## Abstract
We report results of magnetization, zero field resistivity and magnetoresistance measurements in ferromagnetic Heusler alloy Co<sub>2</sub>TiSn. There is a striking change in the character of electron transport as the system undergoes the paramagnetic to ferromagnetic transition. In the paramagnetic state the nature of the electron transport is like that of a semiconductor and this changes abruptly to metallic behaviour at the onset of ferromagnetic ordering. Application of external magnetic field tends to suppress this semiconducting like transport leading to a negative magnetoresistance which reaches a peak in the vicinity of Curie temperature. Comparison is made with the similar unusual behaviour observed in other systems including UNiSn and manganites.
In recent years recognition of half heusler compound NiMnSb as a potential spin-injector material for spintronics applications1 and the discovery of large magnetic shape memory effect in the full Heusler compound Ni<sub>2</sub>MnGa2 have stimulated much research activity on Heusler alloys in general. The full Heusler alloys have composition X<sub>2</sub>YZ forming in L2<sub>1</sub> structure while the half Heusler alloys with composition XYZ forms in C1<sub>b</sub> structure. Here X and Y are transition elements and Z is an sp-element. Traditionally compounds with Mn occupying the Y-site have been drawing the most of the attentions as they form ideal systems for studying localized 3d metallic magnetism, and such studies of course revealed various interesting functionalities in these materials. Non manganese Heusler alloys are interesting too. For example Fe<sub>2</sub>VAl3 and Fe<sub>2</sub>TiSn 4 have been subject of much attention since 1990s due to the possibility of heavy fermion behaviour in them. On the other hand Co-based Heusler alloys Co<sub>2</sub>YZ, where Y=Ti,Zr, Hf etc and Z=Sn or Al, are considered to be good candidates for studying itinerant electron ferromagnetism5 . Of these Co<sub>2</sub>TiSn is particularly interesting because of its similarity with the prototype half metallic system NiMnSb6 . The shape of the spin-polarised total density of states and the dispersion curves for Co<sub>2</sub>TiSn resembles the electronic spectra of half-metallic ferromagnets 7 . Also the magnetic moment per formula unit is close to an integer number (1 $`\mu _B`$), which should be the case for half-metallic systems. Here we present results of magnetization, resistivity and magnetoresistance study in a polycrystalline sample of Co<sub>2</sub>TiSn. The resistivity shows metallic behaviour below the paramagnetic(PM)- ferromagnetic(FM) transition temperature (T$`{}_{Curie}{}^{}`$ 355K). This is consistent with earlier resistivity studies on Co<sub>2</sub>TiSn between 4 and 300K 8 . The nature of the resistivity behaviour, however, changes from metallic to semiconductor-like at T<sub>Curie</sub>. To the best of our knowledge this metal-semiconductor like transition in Co<sub>2</sub>TiSn has not been reported so far. We show that the resistivity around T<sub>Curie</sub> is quite sensitive to applied magnetic field. Anomalous negative coefficient of resistivity above the magnetic ordering temperature has earlier been observed in various (Fe<sub>1-x</sub>M<sub>x</sub>)<sub>3</sub>Ga (M=V and Ti) and (Fe<sub>1-x</sub>N<sub>x</sub>)<sub>3</sub>Si (N=V, Mn, Ti and Cr) based pseudobinary alloys9 ; 10 and UNiSn 11 . We shall compare our results with these earlier experimental findings. Further we discuss our results in the light of a fairly recent theory on the resistivity due to spin-dependent scattering of carriers in ferromagnetic metals with localized spins12 .
Polycrystalline samples of Co<sub>2</sub>TiSn are prepared by argon-arc melting from starting materials of 4N purity. The sample has been annealed in vacuo at 800<sup>0</sup>C for 7 days. The sample is characterized with standard X-ray diffraction study and it consists of a single phase with cubic L2<sub>1</sub> Heusler structure. Magnetization measurements are performed with a commercial SQUID magnetometer (Quantum Design MPMS-5). Resistivity of the sample is measured with the standard four probe technique. A home made cryocooler is used for zero field resistivity measurements in the temperature regime 30K$``$T$``$385K. A separate oven system is employed to extend the temperature range to 450K. A superconducting magnet with field up to 80 kOe is used for magnetoresistance study.
Fig.1 presents the magnetization (M) versus temperature (T) plot for the polycrystalline Co<sub>2</sub>TiSn sample in fields of 100 Oe, 10 kOe and 50 kOe. The T<sub>Curie</sub> is estimated from the inflection point in the M vs T curve (or even better from the peak of dM/dT vs T curve). The estimated T<sub>Curie</sub> from the 100 Oe M vs T data is $``$350K. This is slightly lower than the reported value of 359K13 .However, it is clear from Fig.1 that with the increase in the applied field the ferromagnetism continues to persist in the higher T regime. Consequently the estimated values of T<sub>Curie</sub> are higher. These results clearly show that the magnetic properties of Co<sub>2</sub>TiSn around the PM-FM transition are quite sensitive to the applied field.
Fig.2 presents the resistivity ($`\rho `$) versus T plot in zero applied field. There is a distinct change in the resistivity from metallic to semiconductor-like taking place at 355K. This temperature matches well with the T<sub>Curie</sub> determined from the magnetization measurements, hence we comment that this change in the character of electronic transport is associated with the FM-PM transition. Although the resistivity in the FM regime shows metallic character, magnitude of the resistivity remains fairly high (above 300 $`\mu \mathrm{\Omega }`$-cm) even in the low temperature regime. However, it was possible to fit this metallic resistivity with a standard Bloch-Gruniessen law plus a T<sup>2</sup> term. The Debye temperature estimated from this fitting is 411K which is comparable to that of the isostructural Heulser compound Fe<sub>2</sub>TiSn 4 . The T<sup>2</sup> term can arise due to the spin-flip scattering of charge carriers by magnons in a ferromagnet, and also due to strong interactions in the Fermi liquid. A careful study of resistivity, at least down to helium temperature, is needed to discern amongst the various possible origins. Such a study is in progress now. On the other hand, the resisvity above 355K can be fitted into an expression e$`^{Eg/2k_BT}`$ representing activated behaviour in a semiconductor. Fig.3 presents logarithm of resistivity versus 1/T plot and a band gap of 12.3$`\pm `$1 mev is estimated from the slope of this plot.
Fig. 4 presents the effect of applied magnetic field on the resistivity in and around the PM-FM transition region. The zero field semiconducting like behaviour tends to get suppressed with the increase of the magnetic field. For the sake of clarity we present in Fig.4 only the results with extreme fields of 0 and 80 kOe. Our preliminary study indicates that the semiconducting gap is reduced by 10% with the application of 10 kOe field. However, we need to extend the present temperature range of field dependent resistivity measurement to make a firm comment on the field dependence of the semiconducting gap. In Fig.5 we plot magnetoresistance obtained from the results presented in Fig.4. Magnetoresistance shows a distinct extremum around T<sub>Curie</sub> which suggests a correlation between the spin orientation and the electron scattering. This conjecture is further supported by the observed change in the temperature dependence of magnetization with H in the same T regime (see Fig.1). We have also measured magnetoresistance as a function of field at various fixed temperatures between 50 and 380K. Magnetoresistance becomes negligibly small below 200K and above 370K.
A negative coefficient of resistance has earlier been reported above the Curie temperatures for Ni-doped Co<sub>2</sub>TiSn alloys8 . In that work the resistivity measurement for pure Co<sub>2</sub>TiSn, however, was up to 300K only, hence only metallic behaviour was reported for this parent compound. It is well known from the band structure calculation that the spin down density of states function of Co<sub>2</sub>TiSn varies strongly with energy at the Fermi level (E<sub>F</sub>). As a consequence a slight change in the position of the Fermi level can cause significant variations in the transport properties and this aspect was investigated in details with doping studies8 . However, no definite explanation was provided for the observed negative coefficient of resistance in the Ni-doped Co<sub>2</sub>TiSn.
Such unusual correlation between semiconductor-metal transition and magnetic transition has earlier been observed in UNiSn and this has been a subject of considerable experimental11 and theoretical14 ; 15 attention. UNiSn, like its isostructural compound NiMnSb, is considered to be a half-metallic ferromagnet14 ; 15 ; 16 . From the resistivity measurements an energy gap of 120 mev has been estimated for the semiconducting state of UNiSn. It is conjectured that at the onset of the magnetic transition the energy gap of at least one of the spin bands disappears leading to the metallic behaviour in resistivity. This picture is consistent with the half metallic character of the magnetic state. These results can prompt to think of a one-to-one correspondence between UNiSn and Co<sub>2</sub>TiSn. However, some notes of caution here: first the half metallic nature of Co<sub>2</sub>TiSn is yet to be established firmly, and second there remains some controversy on the magnetic ground state of UNiSn whether it ferromagnetic11 or antiferromagnetic15 ; 16 .
It is interesting to note here that the electronic structure calculations in the sister compound Fe<sub>2</sub>TiSn predicted a nonmagnetic ground state and the existence of pseudogap at the Fermi level4 . In real systems, lattice disorder complicated the possibility of a metal-semiconductor transition; a weak ferromagnetic ground state was observed below 240 K followed by semimetal to semiconductor like transition around 50K 4 . In addition there is evidence of heavy fermion behaviour in Fe<sub>2</sub>TiSn in the low T regime 4 . Effect of Co-doping in Fe<sub>2</sub>TiSn has been studied in details and a negative coefficient of resistivity has been reported for FeCoTiSn for a substantial temperature regime starting from room temperature 17 . However, no definite explanation is available for such behaviour. It is also interesting to note that the Heusler alloy Fe<sub>2</sub>VAl which is also on the verge of a magnetic ordering shows semiconducting behaviour down to 2K 3 . While initial specific heat measurements3 indicated the possibility of electronic mass enhancement in this compound, latter studies suggested that sample-dependent Schottky anomaly originating from magnetic clusters associated with the Fe-defects might be at the origin of the anomalous specific heat behaviour 18 .
A negative coefficient of resistivity above Curie temperature has earlier been reported for ferromagnetic systems like V and Ti-doped Fe<sub>3</sub>Ga 9 , V, Ti, Mn and Cr-doped Fe<sub>3</sub>Si pseudobinary alloys10 . A combination of small conduction electron number per atom and very large spin-disorder scattering was proposed to be the cause of such anomalous resistivity behaviour above the Curie temperature. It was experimentally observed that only those alloys with the resistivity in the paramagnetic state above 150 $`\mu \mathrm{\Omega }`$-cm showed the unusual resistivity behaviour 10 . We note here that the resistivity value of $``$ 500 $`\mu \mathrm{\Omega }`$-cm in the paramagnetic state of Co<sub>2</sub>TiSn fulfills this criterion. A more recent experiment on the transition metal compound MnRhP has revealed all these interesting features 19 .
Changing sign of the coefficient of resistivity at the PM-FM transition is one of the remarkable features of many manganite compounds showing colossal magneto resistance (CMR). The magnitude of the resistivity change and its sensitivity to applied H around T<sub>Curie</sub> is very high in CMR-manganites, hence leading to the drastic CMR effects 20 . While we note that the magnetoresistance in Co<sub>2</sub>TiSn is definitely not comparable with the CMR observed in various manganites, qualitatively the observed peak in magnetotransport around T<sub>Curie</sub> in Co<sub>2</sub>TiSn is very similar in nature. Magnetic interactions and underlying microscopic origin of the magnetic phase transition in Heusler alloys are expected to be quite different from those in manganites and underlying connections between these two different classes of materials are not very obvious. However, there is fairly recent theoretical effort to understand such behaviour in terms of spin-dependent scattering of carriers in ferromagnetic materials with localized spins 12 . In this theoretical framework carrier concentration is a key factor and so is the spin fluctuation especially when a ferromagnetic state is almost unstable against another magnetic state. This theory suggests that in the systems with small Fermi surface, critical spin fluctuations with long wavelengths can contribute to the resistivity to give rise to a peak at T<sub>Curie</sub>. This peak becomes more sharp in a metal whose ferromagnetic state is on the verge of an instability. External magnetic fields can easily suppress critical spin fluctuations with long wavelengths, resulting in a negative magnetoresistance. A ferromagnetic metal with a small Fermi surface can often become half-metallic because of splitting of the up and down-spin bands due to magnetization. While Co<sub>2</sub>TiSn has a low carrier concentration and it is supposedly a half metallic ferromagnet, it is not very clear to the present authors whether detail characteristics of the Fermi surface of Co<sub>2</sub>TiSn would really fit within such theoretical picture.
In conclusion, combining the results of magnetic and transport measurements we show that Co<sub>2</sub>TiSn undergoes a semiconductor-metal transition around 350K and this is associated with a PM-FM transition. In the light of the existing results on UNiSn and recent theoretical studies on spin-dependent scattering of carriers in ferromagnetic materials, the metallic behaviour in the ferromagnetic state of Co<sub>2</sub>TiSn is consistent with the conjectured half-metallic character of this compound. It seems that Co<sub>2</sub>TiSn belongs to a growing class of ferromagnets with low carrier concentration which shows such unusual correlation between magnetic ordering and metal-semiconductor transition.
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# Transverse Momentum Correlations in Relativistic Nuclear Collisions
## 1 Introduction
Central Au-Au collisions at RHIC may generate a color-deconfined medium (quark-gluon plasma or QGP) . Some theoretical descriptions predict abundant low-$`Q^2`$ gluon production in the early stages of high-energy nuclear collisions, with rapid parton thermalization as the source of the colored medium . Nonstatistical fluctuations of event-wise mean $`p_t`$ $`p_t`$ may isolate fragments from low-$`Q^2`$ partons and determine the properties of the corresponding medium. A recent measurement of excess $`p_t`$ fluctuations in Au-Au collisions at 130 GeV revealed a large excess of fluctuations compared to independent-particle $`p_t`$ production .
In this paper we describe the event-wise structure of transverse momentum $`p_t`$ produced in relativistic nuclear collisions at RHIC. We discuss the role of low-$`Q^2`$ partons as Brownian probe particles in heavy ion collisions. We compare joint autocorrelations on $`(\eta ,\varphi )`$ to conventional leading-particle techniques for parton fragment analysis. We present experimental evidence from mean-$`p_t`$ fluctuations and corresponding $`p_t`$ autocorrelations for local temperature/velocity structure in A-A collisions which can be interpreted in terms of parton dissipation in the A-A medium and same-side recoil response of the bulk medium to parton stopping. Finally, we review the energy dependence of mean-$`p_t`$ fluctuations from SPS to RHIC and its implications.
## 2 Low-$`Q^2`$ partons as Brownian probes
In 1905 the microscopic structure of ordinary matter was addressed theoretically by Einstein, who introduced the concept of a (Brownian) probe particle large enough to be observed visually, yet small enough that its motion in response to the molecular dynamics of a fluid might also be observed . Those two constraints specified the one-micron probe particles used by Jean Perrin to confirm molecular motion in fluids . The Langevin equation $`\dot{\stackrel{}{v}}(t)=\frac{1}{\tau }\stackrel{}{v}(t)+\stackrel{}{a}_{stoch}(t)+\stackrel{}{a}_{mcs}(t)`$ models the motion of a Brownian probe in a thermalized fluid medium of point masses qualitatively smaller than the probe particle . The accelerations are gaussian-random with zero mean; $`\stackrel{}{a}_{stoch}(t)`$ is isotropic and $`\stackrel{}{a}_{mcs}(t)\stackrel{}{v}(t)`$ (and $`v`$). The first term models collective dissipation of probe motion (viscosity), the second models individual probe collisions with medium particles and the third simulates multiple Coulomb scattering of a fast probe particle. A solution of that equation for unit initial speed in the $`x`$ direction starting at the ($`x,y`$) origin is shown in the first two panels of Fig. 1. Speed is dissipated with time, leading to equilibration with the medium: fluctuations of velocity about zero and random walk of the probe. An example of such motion is shown in the third panel: an electron track in a time projection chamber exhibits multiple Coulomb scattering along its trajectory, terminating in random walk represented by the circled ball of charge at the trajectory endpoint .
In 2005 we seek the microscopic structure and local dynamics of the QCD medium formed in RHIC heavy ion collisions. The point-mass concept of Einstein’s Brownian probe must be extended to partonic probes, possibly with internal degrees of freedom and experiencing complex non-point interactions with medium degrees of freedom. This problem requires novel analysis techniques closely coupled to the Langevin equation and its associated numerical methods. The analog in heavy ion collisions to Einstein’s Brownian probe is the low-$`Q^2`$ parton, visualized for the first time by methods presented in this paper. In contrast to Einstein’s notion of a particle of exceptional size observed indefinitely in equilibrium with microscopic motions of a thermalized particulate medium, the QCD Brownian probe is identical to medium particles but possesses an exceptional initial velocity relative to the medium with which it interacts for a brief interval: do probe manifestations in the hadronic system reveal ‘microscopic’ degrees of freedom of the medium, is the medium locally or globally equilibrated, what are its fluid properties?
## 3 Joint autocorrelations vs conditional distributions
Conventional study of QCD jets in elementary collisions is inherently model-dependent. Scattered partons with large transverse momentum are associated individually with concentrations of transverse momentum or energy localized on angle variables $`(\eta ,\varphi )`$. In heavy ion collisions, where such identification is impractical, jet studies are based on a high-$`p_t`$ ‘leading particle’ which may estimate a parton momentum direction and some fraction of its magnitude. The leading-particle momentum is the basis for two-particle conditional distributions on transverse momentum and angles. Those distributions reveal medium modifications to parton production and fragmentation as changes in the single-particle $`p_t`$ spectrum ($`R_{AA}`$) and in the fragment-pair relative azimuth distribution (away-side jet disappearance), referred to collectively as jet quenching . The leading-particle approach is based on perturbative concepts of parton hard scattering as a point-like binary interaction and parton energy loss as gluon bremsstrahlung. We can then ask how the medium is modified by parton energy loss and what happens to low-$`Q^2`$ partons, in a $`Q^2`$ regime where the pQCD assumption of point-like interactions breaks down, where the parton may have an effective internal structure. In other words, how can we describe parton dissipation as a transport process, including bulk-medium degrees of freedom?
To access low-$`Q^2`$ partons we have developed an alternative analysis method for jet correlations employing autocorrelation distributions which do not require a leading- or trigger-particle concept. The autocorrelation principle is illustrated in Fig. 2. Projections of the two-particle momentum space of 130 GeV Au-Au collisions onto subspaces $`(\eta _1,\eta _2)`$ and $`(\varphi _1,\varphi _2)`$ (left panels) indicate that correlations on those spaces are approximately invariant on sum variables $`\eta _\mathrm{\Sigma }\eta _1+\eta _2`$ and $`\varphi _\mathrm{\Sigma }\varphi _1+\varphi _2`$, in which case autocorrelations on difference variables $`\eta _\mathrm{\Delta }\eta _1\eta _2`$ and $`\varphi _\mathrm{\Delta }\varphi _1\varphi _2`$ retain nearly all the information in the unprojected distribution . The autocorrelation concept was first introduced to solve the Langevin equation, to extract deterministic information from stochastic trajectories. In time-series analysis the autocorrelation of time series $`f(t)`$ is $`A(\tau )=\frac{1}{T}_{T/2}^{T/2}f(t)f(t+\tau )𝑑t`$, where difference variable $`\tau t_1t_2`$ is the lag. For a stationary distribution ($`f(t)`$ correlations statistically independent of absolute time) the autocorrelation represents a projection by averaging of all the information in $`f(t)`$.
The same principle can be applied to ensemble-averaged two-particle momentum distributions which are approximately invariant on their sum variables . Distributions on angle space $`(\eta _1,\eta _2,\varphi _1,\varphi _2)`$ can be reduced to joint autocorrelations on difference variables $`(\eta _\mathrm{\Delta },\varphi _\mathrm{\Delta })`$. For example, joint autocorrelations in the right-most two panels of Fig. 2 correspond to $`(\eta _1,\eta _2)`$ and $`(\varphi _1,\varphi _2)`$ distributions in the left-most four panels. Joint autocorrelations for relativistic nuclear collisions retain almost all correlation structure on a visualizable 2D space and provide access to parton fragment angular correlations with no leading-particle condition, sampling a minimum-bias parton distribution. Jet correlations are thus revealed with no a priori jet hypothesis, providing access to the low-$`Q^2`$ partons which serve as Brownian probes of the QCD medium.
## 4 The p-p reference system
The reference system for low-$`Q^2`$ partons in A-A collisions is the hard component of correlations in p-p collisions. The single-particle $`p_t`$ spectrum for p-p collisions can be decomposed into soft and hard components on the basis of event multiplicity dependence . Event multiplicity determines statistically the fraction of p-p collisions containing observable parton scattering (hard component). Hard components for ten multiplicity classes in the first panel of Fig. 3, obtained by subtracting fixed soft-component spectrum model $`S_0`$, are plotted on transverse rapidity $`y_t\mathrm{ln}\{(m_t+p_t)/m_0\}`$. The approximately gaussian distributions on $`y_t`$ may be compared with conventional fragmentation functions plotted on logarithmic variable $`\xi \mathrm{ln}\{E_{jet}/p_t\}`$ . Such single-particle structures motivated a study of two-particle correlations on $`(y_{t1},y_{t2})`$. An example in Fig. 3 (second panel) reveals structures at smaller and larger $`y_t`$.
Soft and hard correlation components on $`y_t`$, interpreted as longitudinal string fragments (smaller $`y_t`$) and transverse parton fragments (larger $`y_t`$), produce corresponding structures in joint angular autocorrelations on $`(\eta _\mathrm{\Delta },\varphi _\mathrm{\Delta })`$. In the third panel, string-fragment correlations for unlike-sign pairs are determined by local charge and transverse-momentum conservation (the sharp peak at the origin is conversion electrons). Minimum-bias parton fragments in the fourth panel produce classic jet correlations, with a same-side ($`\eta _\mathrm{\Delta }<\pi /2`$) jet cone at the origin and an away-side ($`\eta _\mathrm{\Delta }>\pi /2`$) ridge corresponding to the broad distribution of parton-pair centers of momentum. Similar-quality parton fragment distributions on $`(\eta _\mathrm{\Delta },\varphi _\mathrm{\Delta })`$ can be obtained for both $`p_t`$s of a hadron pair down to 0.35 GeV/c (parton $`Q/2`$ 1 GeV). The criteria for partons as Brownian probes are 1) $`Q^2`$ large enough that resulting hadron correlations are statistically significant and uniquely assigned to parton fragments, and 2) $`Q^2`$ small enough that correlations are significantly modified by local medium dynamics. In the QCD context the medium itself is formed from low-$`Q^2`$ partons. In the low-$`Q^2`$ regime ‘partons’ may not interact as point color charges, and complex couplings to the medium, e.g., tensor components of the velocity field (Hubble expansion), may be important. Non-perturbative aspects of low-$`Q^2`$ parton collisions should be accessible via low-$`p_t`$ fragment angular autocorrelations and two-particle $`y_t`$ distributions.
## 5 $`p_t`$ fluctuations and prehadronic temperature/velocity structure
Event-wise $`p_t`$ fluctuations generally result from local event-wise changes in the shape of the single-particle $`p_t`$ spectrum, as illustrated in Fig. 4 (first panel). In each collision, a distribution of ‘source’ temperature and/or velocity on $`(\eta ,\varphi )`$ determines the local parent $`p_t`$ spectrum shape. Each hadron $`p_t`$ samples a spectrum shape determined by the sample location, as shown in Fig. 4 (second panel). The local parent shape can be characterized schematically by parameter $`\beta (\eta ,\varphi )`$, interpreted loosely as $`1/T`$ or $`v/c`$ for the local pre-hadronic medium. Variation of either or both parameters relative to an ensemble mean results in $`p_t`$ fluctuations.
A similar situation is encountered in studies of the cosmic microwave background (CMB) as shown in Fig. 4 (third panel) . The temperature distribution $`\beta `$ on the unit sphere is represented by the microwave power density (local spectrum integral rather than mean). The $`\beta (\theta ,\varphi )`$ structure for that single event is directly observable due to large photon numbers. In contrast, for a single heavy ion collision as in Fig. 4 (fourth panel) the parent distribution is sparsely sampled by $`1000`$ final-state hadrons, and parent properties are not accessible on an event-wise basis.
Interpreting $`p_t`$ fluctuations has two aspects: 1) study equivalent two-particle number correlations on $`p_t`$ or $`y_t`$, which reveal medium modification of the two-particle parton fragment distribution—those correlations are directly related to a distribution on $`(\beta _1,\beta _2)`$ sensitive to in-medium parton dissipation; 2) invert the scale or bin-size dependence of $`p_t`$ fluctuations to obtain $`p_t`$ autocorrelations on $`(\eta ,\varphi )`$ which reveal details of event-wise $`\beta (\eta ,\varphi )`$ distribution. We first consider properties of $`\beta (\eta ,\varphi )`$ as a random variable and its relation to two-particle correlations on $`p_t`$ or $`y_t`$. We then employ $`p_t`$ autocorrelations from $`p_t`$ fluctuations to infer aspects of the $`\beta (\eta ,\varphi )`$ distribution which depend only on separation of pairs of points on $`(\eta ,\varphi )`$.
## 6 Parton Dissipation in the A-A Medium
$`p_t`$ fluctuations can be related to a 1D distribution on temperature/velocity parameter $`\beta `$ and corresponding two-point distribution on $`(\beta _1,\beta _2)`$. Each entry of those distributions corresponds to an event-wise $`p_t`$ spectrum in a single bin or pair of bins on $`(\eta ,\varphi )`$. The frequency distribution on $`\beta `$ represents variation of the single-particle $`p_t`$ spectrum shape. For Gaussian-random fluctuations the relative variance of the $`\beta `$ distribution is $`\sigma _\beta ^2/\beta _0^21/n`$, where $`n`$ is the exponent of Lévy distribution $`A/(1+\beta _0(m_tm_0)/n)^n`$ describing the average $`p_t`$ spectrum shape . The shape of the single-particle spectrum is thus related to the event-wise temperature/velocity distribution. Other aspects of shape determination, such as collective radial flow, also contribute to exponent $`n`$. We therefore consider the two-point distribution on $`(\beta _1,\beta _2)`$.
Given the correspondence between the fluctuation distribution on $`\beta `$ and the shape of the single-particle spectrum on $`p_t`$ we seek the relation between the distribution on $`(\beta _1,\beta _2)`$ and the shape of the two-particle distribution on $`(p_{t1},p_{t2})`$. The distribution on $`(\beta _1,\beta _2)`$ provides information about the correlation structure of event-wise $`\beta `$ distributions. The two-particle Lévy distribution on $`(p_{t1},p_{t2})`$, constructed as a Cartesian product of two single-particle distributions with Lévy exponent $`n`$, represents a mixed-pair reference distribution (pairs from different but similar events). We can also define a two-particle object Lévy distribution representing sibling pairs (pairs formed from single events), with exponents $`n_\mathrm{\Sigma }`$ and $`n_\mathrm{\Delta }`$ representing variances on sum and difference axes $`(\beta _\mathrm{\Sigma },\beta _\mathrm{\Delta })`$. The ratio of object and reference distributions reveals a saddle-shaped structure whose curvatures measure temperature/velocity correlations on $`(\eta ,\varphi )`$.
Ratios of sibling to mixed pair densities for 130 GeV Au-Au collisions are shown in Fig. 5 (first two panels) plotted on variable $`X(p_t)`$ . Those panels are dominated by a Lévy saddle, a 2D manifestation of two-particle $`p_t`$ spectrum shape variation due to velocity and temperature fluctuations in the parent distribution. The saddle is an intermediate shape in the dissipation process; its curvatures reflect the correlation structure of the $`(\beta _1,\beta _2)`$ distribution, especially its covariance as discussed in . The saddle curvatures on sum and difference variables, measured by $`1/n_\mathrm{\Sigma }1/n`$ and $`1/n_\mathrm{\Delta }1/n`$, represent the variance excesses (beyond independent $`p_t`$ sampling from a fixed parent) and covariance of temperature/velocity fluctuations for small-amplitude Gaussian-random fluctuations. For an equilibrated system the saddle would be flat (zero curvatures), and $`\beta `$ fluctuations would be consistent with finite-number fluctuations: $`\sigma _{\beta _\mathrm{\Sigma }}^2/\beta _0^2=\sigma _{\beta _\mathrm{\Delta }}^2/\beta _0^2=1/n`$. The integral of correlations on $`(p_{t1},p_{t2})`$, measured by the saddle-curvature difference $`1/n_\mathrm{\Sigma }1/n_\mathrm{\Delta }`$, is equivalent to $`p_t`$ fluctuations measured in the corresponding detector acceptance . With increasing Au-Au centrality the curvature on the difference axis increases strongly, while that on the sum axis approaches zero .
More recently, we have transitioned from per-pair correlation measure $`\widehat{r}1`$ plotted on variable $`X(p_t)`$ to per-particle density ratio $`\mathrm{\Delta }\rho /\sqrt{\rho _{ref}}`$ plotted on transverse rapidity $`y_t`$. We wish to follow, within a single context, the transition from parton fragment distributions in elementary collisions to correlations from parton dissipation in a bulk medium. Fig. 4 (last two panels) shows $`\mathrm{\Delta }\rho /\sqrt{\rho _{ref}}`$ on $`(y_{t1},y_{t2})`$ for peripheral and central Au-Au collisions at 200 GeV. The logarithmic $`y_t`$ interval \[1,4.5\] corresponds to linear $`p_t`$ \[0.15,6\] GeV/c. Peripheral collisions produce a 2D minimum-bias parton fragment distribution peaked at $`y_t2.5`$ ($`p_t1`$ GeV/c), similar to p-p collisions but without small-$`y_t`$ correlations from string fragmentation. As centrality increases the fragment distribution is transported to smaller $`y_t`$ and approaches a shape corresponding to the Lévy saddle on $`X(p_t)\times X(p_t)`$. In this format we can study the transition with A-A centrality between two extreme cases: 1) in vacuo distributions of string and parton fragments and 2) gaussian-random variation of $`\beta `$ on $`(\eta ,\varphi )`$ for a nearly-equilibrated system. Parton dissipation in the A-A bulk medium is represented by the transition between those extremes.
## 7 $`p_t`$ fluctuations and $`p_t`$ autocorrelations
The previous section describes $`p_t`$ fluctuations in terms of two-particle number densities on $`(p_{t1},p_{t2})`$ or its logarithmic equivalent $`(y_{t1},y_{t2})`$, the issue being modification of the two-particle parton fragment distribution with changing A-A centrality. One can also express $`p_t`$ fluctuations in terms of two-particle $`p_t`$ distributions on $`(\eta ,\varphi )`$ which reveal different aspects of the underlying two-particle number distribution on vector momentum. This section describes a procedure to determine the correlation structure of the $`\beta (\eta ,\varphi )`$ distribution as a temperature/velocity distribution on the prehadronic medium.
Fluctuations in bins of a given size or scale are determined by two-particle correlations with characteristic lengths less than or equal to the bin scale. By measuring fluctuation magnitudes as a function of bin size one can recover some details of the two-particle correlation structure—those aspects which depend on the separation of pairs of points, not on their absolute positions. The relation between fluctuations and correlations is given by the integral equation
$`\mathrm{\Delta }\sigma _{p_t:n}^2(mϵ_\eta ,nϵ_\varphi )=4{\displaystyle \underset{k,l=1}{\overset{m,n}{}}}ϵ_\eta ϵ_\varphi `$ $`K_{mn;kl}{\displaystyle \frac{\mathrm{\Delta }\rho (p_t:n;kϵ_\eta ,lϵ_\varphi )}{\sqrt{\rho _{ref}(n;kϵ_\eta ,lϵ_\varphi )}}},`$ (1)
with kernel $`K_{mn;kl}(mk+1/2)/m(nl+1/2)/n`$ representing the 2D macrobin system, $`\mathrm{\Delta }\sigma _{p_t:n}^2(\delta \eta ,\delta \varphi )`$ is a variance excess and $`\mathrm{\Delta }\rho (p_t:n;)/\sqrt{\rho _{ref}(n)}`$ is an autocorrelation density ratio. That equation can be inverted numerically to obtain the $`p_t`$ autocorrelation.
Fig. 6 shows fluctuation scale dependence on bin sizes $`(\delta \eta ,\delta \varphi )`$ and joint $`p_t`$ autocorrelations on difference variables $`(\eta _\mathrm{\Delta },\varphi _\mathrm{\Delta })`$ for peripheral (left panels) and mid-central (right panels) Au-Au collisions. Fluctuation measurements at the full STAR acceptance correspond to the single points at the apex of the distributions on scale. Measurements with different detectors correspond to different regions of those surfaces. Inversion to autocorrelations provides physical interpretation of fluctuation scale dependence. By inverting $`p_t`$ fluctuations, parton fragment distributions are visualized as temperature/velocity structures on $`(\eta ,\varphi )`$ complementary to number correlations on $`(y_{t1},y_{t2})`$ described in the previous section. A more comprehensive picture of parton scattering, dissipation and fragmentation is thereby established.
$`p_t`$ autocorrelations can also be determined directly by pair counting. In Fig. 7 the peripheral Au-Au result from the previous section (first panel) is compared to the minimum-bias p-p result (second panel) and to p-p collisions with $`n_{ch}9`$ (third panel). The last panel shows the charge-dependent (like-sign $``$ unlike-sign pairs) $`p_t`$ autocorrelation for the same event class, reflecting charge-ordering along the jet thrust axis during parton fragmentation. This is the first determination of $`p_t`$ correlations in p-p collisions.
## 8 Local velocity structure and same-side recoil
Whether derived from pair counting or from fluctuation inversion, the resulting $`p_t`$ autocorrelations can be separated into several components. We first subtract multipoles on azimuth (azimuth sinusoids independent of pseudorapidity), revealing structure associated with parton scattering and fragmentation. Fig. 8 shows the resulting $`p_t`$ autocorrelation for 20-30% central Au-Au collisions at 200 GeV (first panel) and a three-component model fit to that distribution (second panel) including a same-side ($`\varphi _\mathrm{\Delta }<\pi /2`$) positive peak, a same-side negative peak and an away-side ($`\varphi _\mathrm{\Delta }>\pi /2`$) positive peak. The fit is excellent, with residuals at the percent level. The third panel shows the result of subtracting the positive same-side model peak (representing parton fragments) from the data in the first panel. The shape of the negative same-side peak is very different from the positive peak; there is thus negligible systematic coupling in the fit procedure. The fourth panel shows the data distribution in the third panel plotted in a cylindrical format, suggesting an interpretation in terms of temperature/velocity correlations.
Histogram values of the $`p_t`$ autocorrelation effectively measure correlations (covariances) of blue or red shifts of local $`p_t`$ spectra relative to the ensemble mean spectrum at pairs of points separated by ($`\eta _\mathrm{\Delta },\varphi _\mathrm{\Delta }`$). The negative same-side peak can therefore be interpreted as a systematic red shift of local $`p_t`$ distributions adjacent to the positive fragment peak. The red shift can in turn be interpreted as recoil of the bulk medium in response to stopping the parton partner of the observed parton (positive same-side peak). This detailed picture of parton dissipation, stopping and fragmentation in a A-A collisions, including recoil response of the dissipative bulk medium suggested in the fourth panel, is accessed for the first time with joint $`p_t`$ autocorrelations.
## 9 Reconstructing Event-wise Temperature/velocity Structure
We now consider the relation of $`p_t`$ autocorrelations to individual collision events. In Fig. 9 we repeat the WMAP CMB distribution of microwave power on the unit sphere, picturing a single Big Bang ‘event’ which has a large statistical depth and can therefore be directly observed . Information relevant to cosmological theory is extracted as a power spectrum on polar angle (second panel), formally equivalent (within a Fourier transform) to an autocorrelation according to the Wiener-Khinchine theorem. In some studies, CMB angular autocorrelations and cross-correlation have been determined directly .
In our study of heavy ion collisions we obtain angular $`p_t`$ autocorrelations as in the third panel. Due to sparse sampling we cannot directly visualize the temperature/velocity structure of individual collision events as for the CMB survey. The local microwave power density of the CMB survey is analogous to local $`p_t`$ in a Au-Au collision. For individual collisions, and especially for smaller bin sizes, the event-wise mean values are not significant. However, given $`p_t`$ autocorrelations we can simulate event-wise velocity/temperature distributions. We estimate the number of hard parton scatters within the STAR acceptance in a central Au-Au collision as 20-40, based on an analysis of p-p collisions . Combining that frequency estimate with shape information from the autocorrelation, and introducing some statistical variation of peak structure about the autocorrelation mean value, we can produce simulated events as shown in Fig. 9 (fourth panel): distributions on primary angle variables $`(\eta ,\varphi )`$, whereas the autocorrelation is on difference variables $`(\eta _\mathrm{\Delta },\varphi _\mathrm{\Delta })`$. This exercise illustrates that while Au-Au collisions are RHIC may be locally equilibrated prior to kinetic decoupling, they remain highly structured due to copious parton scattering which is not fully dissipated. Access to that structure requires $`p_t`$ rather than angular or number autocorrelations on $`(\eta ,\varphi )`$ to provide the full picture.
## 10 Energy dependence of $`p_t`$ fluctuations and parton scattering
Given this close connection between parton scattering and fluctuations, the collision-energy dependence of $`p_t`$ fluctuations may reveal previously inaccessible parton dynamics at lower collision energies. In Fig. 10 (first panel) we show the centrality dependence ($`\nu `$ measures mean participant path length in nucleon diameters) of $`p_t`$ fluctuations for four RHIC energies and a summary (crosshatched region) of SPS fluctuation measurements at 12.6 and 17.3 GeV , all at the full STAR acceptance (CERES measurements are extrapolated). In the second panel the pseudorapidity scale dependence of fluctuations at full azimuth acceptance is shown for central collisions at six energies. Extrapolation of CERES data in the first panel is illustrated by the dashed lines at the bottom of the second. Fluctuation measure $`\mathrm{\Delta }\sigma _{p_t:n}`$ is related to the variance difference in Eq. (1) by $`\mathrm{\Delta }\sigma _{p_t:n}^22\sigma _{\widehat{p}_t}\mathrm{\Delta }\sigma _{p_t:n}`$, with $`\sigma _{\widehat{p}_t}`$ the single-particle variance. To good approximation $`\mathrm{\Delta }\sigma _{p_t:n}\mathrm{\Phi }_{p_t}`$, and both are per particle fluctuation measures. $`\mathrm{\Phi }_{p_t}`$ was used for the CERES fluctuation measurements.
For either measure we observe a dramatic increase in $`p_t`$ fluctuations from SPS to RHIC energies. The centrality dependence in the first panel suggests that fluctuations for p-p and peripheral A-A collisions saturate near 60 GeV, whereas there is monotonic increase for the more central collisions. The scale dependence in the second panel illustrates how measurements with different detector acceptances are related. Measurements over common scale intervals should correspond. At RHIC energies we have demonstrated that $`p_t`$ fluctuations are dominated by fragments from low-$`Q^2`$ parton collisions. The energy dependence of $`\mathrm{\Delta }\sigma _{p_t:n}`$ or $`\mathrm{\Phi }_{p_t}`$ is shown in the third panel of Fig. 10, plotted vs $`\sqrt{s_{NN}}`$. We observe that $`p_t`$ fluctuations vary almost linearly with $`\mathrm{log}\{\sqrt{s_{NN}}/10.5\}`$ (solid curve in that panel), suggesting a threshold for observable parton scattering and fragmentation near 10 GeV.
Fluctuation measurements based on $`\mathrm{\Sigma }_{p_t}\sqrt{\mathrm{\Delta }\sigma _{p_t:n}^2/(\overline{n}_{ch}\widehat{p}_t^2)}`$ appear to contradict the results described here, implying instead negligible energy dependence of $`p_t`$ fluctuations from SPS to RHIC. We observe that nuclear collisions at RHIC are dominated by local temperature/velocity structure from hard parton scattering. $`\mathrm{\Sigma }_{p_t}`$ is a per pair measure which averages the local $`p_t`$ correlation structure dominating RHIC collisions over the entire detector acceptance, resulting in apparent reduction of correlations with increasing A-A centrality as $`1/N_{participant}`$ (per the central limit theorem) and consequent insensitivity to contributions from hard scattering. We want to study separately the changes in $`p_t`$ production ($`T`$) and in the correlation structure of that produced $`p_t`$ ($`\delta T`$) prior to hadronization. $`\mathrm{\Sigma }_{p_t}`$ by construction estimates relative temperature fluctuations of the form $`\delta T/T`$. It thus divides the structure problem by the production problem, greatly decreasing sensitivity to each.
## 11 Summary
We have demonstrated that low-$`Q^2`$ partons, accessed here for the first time by novel analysis techniques including joint autocorrelations, serve as Brownian probes of A-A collisions, being the softest detectable dynamical objects which experience QCD interactions as color charges. Our analysis of p-p correlations provides an essential reference for A-A collisions. Inversion of the scale dependence of $`p_t`$ fluctuations provides the first access to $`p_t`$ autocorrelations which reveal a complex parton dissipation process in A-A collisions relative to p-p collisions. We observe possible evidence for bulk-medium recoil in response to parton stopping. We also observe strong energy dependence of $`p_t`$ fluctuations, which is to be expected given the dominant role of scattered partons in driving those fluctuations.
This work was supported in part by the Office of Science of the U.S. DoE under grant DE-FG03-97ER41020.
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# How an antenna launches its input power into radiation: the pattern of the Poynting vector at and near an antenna
## Appendix A SURFACE ELECTRIC FIELD AS EXPANSION IN ASSOCIATED LEGENDRE FUNCTIONS
A spherical antenna has a gap in its perfectly conducting surface defined by $`\mathrm{cos}\theta _1<\mathrm{cos}\theta <\mathrm{cos}\theta _2`$. The internal source of power creates an electric field $`E_\theta `$ at $`r=a`$ within the gap, uniformly in azimuth. Otherwise, $`E_\theta =0`$ on the surface. The multipole expansion (16) of $`E_\theta `$ is in terms of the associated Legendre functions $`P_{\mathrm{}}^1`$. We thus require an expansion in those associated Legendre functions of the rectangular function, $`f(z)=[\mathrm{\Theta }(zz_1)\mathrm{\Theta }(zz_2)]`$ where $`\mathrm{\Theta }(x)`$ is the Heaviside step function. We begin with the completeness relation on the interval $`(1,1)`$ in $`z=\mathrm{cos}\theta `$ for the Legendre polynomials $`P_{\mathrm{}}(z)`$ and $`P_{\mathrm{}}^1(z)=\sqrt{1z^2}dP_{\mathrm{}}(z)/dz`$:
$`\delta (zz^{})`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\mathrm{}=0}{\overset{\mathrm{}}{}}}(2\mathrm{}+1)P_{\mathrm{}}(z^{})P_{\mathrm{}}(z)`$ (36)
$`\delta (zz^{})`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\mathrm{}=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(2\mathrm{}+1)}{\mathrm{}(\mathrm{}+1)}}P_{\mathrm{}}^1(z^{})P_{\mathrm{}}^1(z)`$ (37)
We integrate (36) in *z* over the interval $`(1,z)`$ to obtain
$$\mathrm{\Theta }(zz^{})=\frac{1}{2}\underset{\mathrm{}=0}{\overset{\mathrm{}}{}}(2\mathrm{}+1)P_{\mathrm{}}(z^{})_1^zP_{\mathrm{}}(z)𝑑z$$
(38)
For $`\mathrm{}=0`$ the integral is
$$\frac{1}{2}_1^zP_0(z)𝑑z=(1+z)/2$$
and for $`\mathrm{}>0`$,
$$\frac{1}{2}_1^zP_{\mathrm{}}𝑑z=\frac{1}{2}_z^1P_{\mathrm{}}𝑑z=\frac{\sqrt{1z^2}}{2\mathrm{}(\mathrm{}+1)}P_{\mathrm{}}^1(z)$$
The final result can be found in Magnus, Oberhettinger, and Soni. MOS Substituting these two results into (38) yields
$$\mathrm{\Theta }(zz^{})=\frac{1+z}{2}+\sqrt{1z^2}\underset{\mathrm{}=1}{\overset{\mathrm{}}{}}\frac{(2\mathrm{}+1)}{2\mathrm{}(\mathrm{}+1)}P_{\mathrm{}}(z^{})P_{\mathrm{}}^1(z)$$
(39)
We define the tangential field on the surface as
$$E_\theta (r=a,z,z_1,z_2)=\frac{A}{\sqrt{1z^2}}[\mathrm{\Theta }(zz_1)\mathrm{\Theta }(zz_2)]$$
(40)
where *A* will be chosen for convenience below. With the expansion (39) for $`\mathrm{\Theta }(zz^{})`$ the tangential electric field in the gap is given by
$$E_\theta (r=a,z,z_1,z_2)=A\underset{\mathrm{}=1}{\overset{\mathrm{}}{}}\frac{(2\mathrm{}+1)}{2\mathrm{}(\mathrm{}+1)}[P_{\mathrm{}}(z_1)P_{\mathrm{}}(z_2)]P_{\mathrm{}}^1(z)$$
(41)
In our calculations we chose the gap to be relatively small and centered around $`\mathrm{cos}\theta =0`$. With $`z_1=\mathrm{cos}(\pi /2+ϵ)`$ and $`z_2=\mathrm{cos}(\pi /2ϵ)`$ we find, using (40), that the voltage *V*, defined as the integral of the electric field across the gap, is $`V=2aϵA`$. From the symmetry of the Legendre functions around $`\theta =\pi /2`$, a symmetric equatorial gap implies only odd $`\mathrm{}`$ terms in (41). The result for the tangential ($`\theta `$-component) electric field on the surface of the sphere and in the gap is
$$E_\theta (r=a,\mathrm{cos}\theta )=\frac{V}{2aϵ}\underset{\mathrm{}odd}{}\frac{(2\mathrm{}+1)}{\mathrm{}(\mathrm{}+1)}P_{\mathrm{}}(\mathrm{sin}ϵ)P_{\mathrm{}}^1(\mathrm{cos}\theta )$$
(42)
In the limit of $`ϵ0`$, equation (42) can be shown to be equal to *V/a* times (37) with $`z^{}=0`$.
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# The Millennium Galaxy Catalogue: Dynamically close pairs of galaxies and the global merger rate
## 1 Introduction
In the Cold Dark Matter (CDM) framework, galaxies are assembled gradually via a process of hierararchical mergers, where increasingly more massive subunits coalesce to produce today’s luminous giant galaxies (e.g., Cole et al. 2000 and references therein): the history of mass assembly of galaxies is then reflected in the global merger rate and in its evolution with lookback time. Mergers and close interactions may also play a role in the morphological transformation of galaxies (see, for instance, Hernandez-Toledo et al. 2005; Scannapieco and Tissera 2003; Steinmetz and Navarro 2002; Junqueira et al. 1998; Mihos and Hernquist 1996), triggering of starbursts (Nikolic et al., 2004; Alonso et al., 2004; Barton et al., 2003; Bergvall et al., 2003; Lambas et al., 2003; Tissera et al., 2002; Barton et al., 2000; Donzelli and Pastoriza, 1997) and fuelling of active galactic nuclei (e.g., Sanchez and González-Serrano 2003; Canalizo and Stockton 2001 ).
Traditionally, the observational route to measuring the merger rate has been the conventional pair fraction, under the assumption that sufficiently close galaxies will result in a merger over relatively short timescales (Zepf and Koo, 1989; Burkey et al., 1994; Carlberg et al., 1994; Woods et al., 1995; Yee and Ellingson, 1995; Patton et al., 1997; Wu and Keel, 1998; Bundy et al., 2004); however, limited redshift information makes the derived pair fraction dependent on the mean galaxy density and the correlation function, and all such estimates are affected by contamination from unphysical pairs.
In order to address this issue, it is sometimes required that pairs also show evidence of interactions (Neuschafer et al., 1997; Le Fevre et al., 2000; Hashimoto and Oemler, 2000; Xu et al., 2004): on the other hand, this introduces an element of subjectivity in the analysis, as a threshold of morphological disturbance must be chosen for objects to be considered parts of pairs. Furthermore, a number of true pairs may not show evidence of tidal tails or other deviations from symmetry (because of low surface brightness, cosmological dimming and morphological $`k`$-corrections, for example; see discussion in Mihos 1995).
The fraction of galaxies in kinematic pairs (i.e. both spatially and dynamically close) yields a more rigorous estimate of the local pair fraction and global merger rate, although the method requires highly complete redshift information and accurate control of systematic biases arising from the flux-limited nature of redshift surveys and from incompleteness. At least for the local universe, where large redshift samples are currently available, this approach is now feasible and was first used by Patton et al. (2000, hereafter P00) who measured the number of dynamically close companions per galaxy within the range $`21<M_B<18`$ (a statistic akin to the pair fraction) using data from the Second Southern Sky Redshift Survey (da Costa et al., 1998), and established the necessary mathematical formalism (see next section for a brief summary).
The two large redshift surveys now available, the 2dF Galaxy Redshift Survey (2dFGRS; Colless et al. 2001), and the Sloan Digital Sky Survey (SDSS; York et al. 2000), should provide a more extensive sample for such studies. Of these, the 2dFGRS is currently the largest and most complete publicly available dataset. However, because of its complicated selection function at small angular separation, which is related to restrictions on fiber placement and the survey tiling strategy, this survey is not well suited to a study of close pairs (Hawkins et al., 2003). In this paper we choose to study a smaller, but more complete dataset, which is more appropriate for the purpose of measuring the number of close companions per galaxy and the local merger rate.
The Millennium Galaxy Catalogue, hereafter MGC (Liske et al., 2003), covers a $`35^{}`$ wide and $`72.2^{}`$ long equatorial strip, coinciding with the northern strip of the 2dFGRS and with the SDSS Data Release 1 region (Abazajian et al., 2003), which provide $`50\%`$ of the MGC redshifts as well as $`ugriz`$ photometry. The MGC has been reimaged in the $`B`$ band to deep surface brightness limits (26 mag arcsec<sup>-2</sup>) in order to derive accurate structural parameters, and additional redshifts have been collected to bring the redshift completeness to 99.79% for galaxies with $`B<19.2`$ and 96.05% for galaxies with $`B<20`$. The total sample includes 10,095 galaxies to $`B=20`$ (Driver et al., 2005). Because of its high completeness and photometric accuracy (Cross et al., 2004), the MGC is well suited to a study of the local merger rate via close pairs and to an analysis of the properties of the dynamically paired galaxies. Even the MGC, however, suffers from some incompleteness at small separations, for which we will need to make allowance in our analysis.
We describe the procedures used to derive the number of close companions per galaxy in the next section, apply these techniques to the MGC in section 3 and present the results and our discussion in section 4, dealing with theoretical interpretation of the data and with the properties of galaxies in pairs. We adopt a cosmology with $`\mathrm{\Omega }_M=0.3`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$ and H$`{}_{0}{}^{}=100`$ km s<sup>-1</sup> Mpc<sup>-1</sup> and all absolute magnitudes used should be understood with reference to this cosmology.
## 2 Methodology
We follow the method of P00 and the reader is referred to that paper (and Patton et al. 2002, hereafter P02) for a more complete discussion. Here we present a summary of the technique and repeat some essential points.
Consider a primary sample of $`N_1`$ galaxies and a secondary sample of $`N_2`$ galaxies. We are interested in determining the number of galaxies from the secondary sample that are dynamically ’close’ to a primary galaxy, per primary galaxy. This statistic, $`N_c`$, is similar to the fraction of galaxies in close pairs and is in fact identical to it for a volume-limited sample which contains only close pairs but no triplets or higher order $`N`$-tuples. We define galaxies to be dynamically close if they have a projected separation $`5<r_p<20`$ kpc (the inner limit is set to avoid confusion with star formation knots or HII regions) and if they have a relative velocity of less than 500 km s<sup>-1</sup>. Simulations show that such pairs would result in a merger on timescales of 10<sup>9</sup> years (Toomre, 1977; Barnes, 1988) and therefore $`N_c`$ is directly related to the galaxy merger rate. These requirements have been conventionally used in the literature (e.g., Carlberg et al. 1994; Yee and Ellingson 1995; Carlberg et al. 2000) on close pair statistics.
P00 have shown that for a flux-limited sample (like the MGC) the number of close companions per galaxy is best computed as:
$$N_c=\frac{\underset{i}{\overset{N_1}{}}w_{N_1}^iN_{c_i}}{_i^{N_1}w_{N_1}^i}$$
(1)
and the total companion luminosity as:
$$L_c=\frac{\underset{i}{\overset{N_1}{}}w_{L_1}^iL_{c_i}}{_i^{N_1}w_{L_1}^i}.$$
(2)
$`N_{c_i}`$ and $`L_{c_i}`$ are the number and luminosity of galaxies from the secondary sample that are dynamically close to the $`i`$th primary galaxy. They are given by:
$$N_{c_i}=\underset{j}{}w_{N_2}^j=\underset{j}{}\frac{w_{b_2}^jw_{v_2}^j}{S_N(z_j)}$$
(3)
and
$$L_{c_i}=\underset{j}{}w_{L_2}^jL_j=\underset{j}{}\frac{w_{b_2}^jw_{v_2}^j}{S_L(z_j)}L_j,$$
(4)
where the sums run over those secondary galaxies that fulfill the criteria of being dynamically close to the $`i`$th primary galaxy.
The weights that are being applied to the secondary sample, $`w_{N_2}`$ and $`w_{L_2}`$, have three components. The first factor corrects for the change of density of the secondary galaxies as a function of redshift due to the apparent flux limit of the sample: for a low-redshift primary galaxy we will find more close companions on average than for a higher redshift one because the companions can be drawn from a wider range of luminosities (cf. Figure 1). The appropriate (inverse) correction factors are given by the selection functions $`S_N(z)`$ and $`S_L(z)`$ which are defined in terms of ratios of densities in flux limited vs. volume limited samples:
$$S_N(z)=\frac{_{M_{bright}}^{M_{lim}(z)}\mathrm{\Phi }(M)𝑑M}{_{M_{bright}}^{M_2}\mathrm{\Phi }(M)𝑑M}$$
(5)
and
$$S_L(z)=\frac{_{M_{bright}}^{M_{lim}(z)}\mathrm{\Phi }(M)L(M)𝑑M}{_{M_{bright}}^{M_2}\mathrm{\Phi }(M)L(M)𝑑M}.$$
(6)
$`\mathrm{\Phi }(M)`$ denotes the luminosity function for which we adopt the latest MGC values from Driver et al. (2005) ($`\mathrm{\Phi }^{}=0.0177\mathrm{Mpc}^3`$, $`M^{}=19.60`$ and $`\alpha =1.13`$; $`B`$ band). $`M_{lim}`$ is the redshift dependent absolute magnitude limit:
$$M_{lim}(z)=\mathrm{max}(M_{faint},m5\mathrm{log}d_L(z)25k(z)e(z)),$$
(7)
where $`m`$ is the apparent magnitude limit of the survey ($`B=20`$), $`d_L`$ is the luminosity distance, $`k(z)`$ is the maximum $`k`$-correction at redshift $`z`$ (P02) and $`e(z)`$ is of the form $`0.75\times 2.5\mathrm{log}(1+z)`$, which corresponds to a passive luminosity evolution scenario with no mergers: this form yields the same results (to within 0.01) of the $`e(z)=0.7z`$ assumed by P02. Here, $`M_{bright}`$ and $`M_{faint}`$ are magnitude limits introduced to take account of the fact that the clustering properties of galaxies vary with luminosity (Norberg et al., 2002) and therefore we must only use galaxies in a limited luminosity range so that the clustering length does not vary too strongly.
In essence, the factors $`S_N`$ and $`S_L`$ in equations 3 and 4 correct the secondary sample from a flux limited sample in the range $`M_{bright}<M<M_{lim}(z)`$ to a hypothetical volume limited sample in the range $`M_{bright}<M<M_2`$.
The other two components of the $`w_{N_2}`$ weights correct for boundary effects. $`w_{b_2}`$ corrects for the fact that some fraction, $`1f_b^i`$, of the $`r_p`$ area around the $`i`$th primary galaxy may lie outside the effective survey area: $`w_{b_2}^j=1/f_b^i`$. For primary galaxies lying close to the survey boundaries in redshift space (i.e. within 500 km s<sup>-1</sup>) $`w_{v_2}`$ corrects for possible companions beyond those limits: we ignore all companions between the primary galaxy and the redshift boundary and adopt $`w_{v_2}^j=2`$ for those companions lying away from the boundary.
Finally we note that in equations 1 and 2 we also apply weights to the primary sample. The purpose of the $`w_{N_1}`$ and $`w_{L_1}`$ weights is to minimise the errors on $`N_c`$ and $`L_c`$. P00 have shown that they should be chosen as:
$$w_{N_1}^i=w_{b_1}^iw_{v_1}^iS_N(z_i)$$
(8)
$$w_{L_1}^i=w_{b_1}^iw_{v_1}^iS_L(z_i),$$
(9)
where $`w_{b_1}^i=f_b^i`$ and $`w_{v_1}^i=0.5`$ (i.e. the reciprocals of the weights for the secondary sample) for primary galaxies close to the redshift boundaries.
## 3 Application to the MGC
We now apply these techniques to the MGC. Figure 1 shows the absolute magnitude vs. redshift for all MGC galaxies with $`B<20`$: this includes 9696 galaxies with redshifts out of 10,095, with 96.05% redshift completeness.
The mean absolute magnitude of the MGC, weighted according to equations (8) and (9) above, is $`M_B=18.5`$. We then choose to analyze two subsets of the MGC data, one with $`M_{bright}=22`$ and $`M_2=19`$ and another with $`M_{bright}=21`$ and $`M_2=18`$ (to which ranges the results are normalized). Selection lines for both samples are drawn in Figure 1. For these samples, we use $`M_{faint}=18`$ and $`17`$ (respectively); the brighter sample contains 5756 galaxies, while the fainter one includes 6492 galaxies. The magnitude limits chosen here (as in other studies) imply that we are mostly concerned with major mergers between galaxies of approximately similar luminosity. As in P02 some galaxies are brighter than the apparent magnitude limit but blue and therefore lie ‘below’ the selection line in Figure 1; these galaxies are excluded from the analysis to ensure that galaxies of all spectral types have an equal probability of being included within the samples.
For the bright sample ($`22<M_B<19`$) we find a total of 137 dynamically close companions in 69 pairs<sup>1</sup><sup>1</sup>1One of the pairs contributes only a single companion because $`r_p`$ is just below 20 kpc at the redshift of one of the galaxies but just above 20 kpc at the redshift of the other, where two galaxies appear in two pairs. Of these companions, 60 lie in the magnitude range $`19<M_B<18`$ (i.e. between $`M_2`$ and $`M_{faint}`$ for this sample. For the faint sample ($`21<M_B<18`$) we find 176 companions in 89 pairs, with one triple and where two galaxies appear in two pairs each. Of these, 27 lie between $`M_2`$ and $`M_{faint}`$. This is similar to P00 in that galaxies between these two limits do not contribute significantly to the pair fraction for the $`L^{}`$ galaxies; however, these objects are much more significant contributors for the brighter sample. This may be evidence that the brighter galaxies tend to have fainter companions.
Applying equations 1 and 2 we find $`N_c=0.0147\pm 0.0013`$ and $`L_c=213\pm 25\times 10^6L_{}`$ for galaxies with $`22<M_B<19`$, with mean $`z=0.124`$ and $`N_c=0.0301\pm 0.0023`$, with $`L_c=(248\pm 26)\times 10^6L_{}`$ for galaxies with $`21<M_B<18`$ with mean $`z=0.116`$, where the errors are derived by jackknife resampling of the data (Efron, 1982).
Figure 2 shows postage stamps for a random selection of pairs; about half or more of these galaxies show clear signs of interaction such as disturbed morphologies, rings and tidal features, demonstrating that our sample indeed includes bona fide physical pairs. The full list of companions provides a useful objectively selected sample of galaxies in the early phase of a merger. Table 1 lists all unique pairs found in MGC; in order the columns show: the MGC ID of the primary galaxy, the IDs of the other pair members, the RA and Dec. (2000) for the primary galaxy and its radial velocity.
However, we know that the MGC has some redshift incompleteness. We proceed to estimate the number of missed pairs in the following way. We repeat our analysis including galaxies without redshifts and using only the $`r_p`$ criterion, but requiring that if a galaxy without redshift is involved in a pair, its absolute magnitude (assigning it the redshift of the other object) be consistent with the absolute magnitude limits being used (i.e. it must lie between the selection lines shown in Figure 1). This yields an extra 36 companions for the $`22<M_B<19`$ sample and 47 companions for the $`21<M_B<18`$ sample. Since some of these pairs may be due to chance superposition, we repeat our analysis using only galaxies with redshifts, but solely using the $`r_p`$ criterion. We find a total of 193 and 262 companions, respectively, of which we know 137 and 176 to be ’real’. The fraction of ‘true’ companions is therefore 70% and we apply this to derive an extra 25.2 companions for the first sample and 32.9 companions for the second sample. This yields a correction of 18.4% and 18.7% to the derived values of $`N_c`$ and $`L_c`$ (for the brighter and fainter sample, respectively) owing to galaxies whose redshifts are missing from the survey.
These values are larger than one might expect from the low redshift incompleteness ($`4\%`$) of the survey; the main source of incompleteness lies in galaxies that were never targeted and these objects lie preferentially at close angular separations to other galaxies: since the MGC derives most of its redshifts from 2dF observations (including 2dFGRS data) and SDSS surveys, the corrections are affected by the instrumental bias against close pairs due to fiber placement constraints (although this is lower than in the 2dFGRS because of the higher MGC completeness).
## 4 Discussion
After applying the incompleteness correction derived above we find that $`N_c=0.0174\pm 0.0015`$ and $`L_c=(252\pm 30)\times 10^6`$ $`L_{}`$ for galaxies with $`22<M_B<19`$ and $`N_c=0.0357\pm 0.0027`$ and $`L_c=(294\pm 31)\times 10^6L_{}`$ for galaxies with $`21<M_B<18`$, with mean redshifts as given above. If we assume that the merger rate evolves as $`(1+z)^m`$ with $`m3`$ (Le Fevre et al., 2000; Gottlöber et al., 2001; Patton et al., 2002) we find that our $`N_c`$ of $`0.0357\pm 0.0027`$ translates to $`0.0257\pm 0.0019`$ at the mean SSRS2 redshift of 0.015, which is in good agreement with the number of companions of $`0.0225\pm 0.0052`$ measured by P00 for galaxies with $`21<M_B<18`$. At $`z=0.289`$, P02 find $`N_c=0.0334\pm 0.0081`$ for these same galaxies, converted to our cosmology (D. Patton, private communication), which corresponds to $`N_c=0.0217\pm 0.053`$ at our mean $`z=0.116`$.
### 4.1 Comparison with theoretical models
Predictions of the expected merger rate in CDM models can be derived from numerical simulations. However, these simulations are usually conducted at much lower resolution than the typical pair separation used here. Although it is possible to extrapolate to higher resolution, by resampling the larger simulations, this procedure is suspicious as information may already have been lost on small scales (Gottlöber et al., 2001). Additionally, it is difficult to relate dark haloes to their luminous content: it is generally assumed that each galaxy corresponds to a single halo, but this assumption is likely to be too simplistic.
Khochfar and Burkert (2001) calculate the fraction of close pairs and the evolution of the merger rate in a $`\mathrm{\Lambda }`$CDM cosmology. The merger fraction they derive for galaxies of $`M3\times 10^{12}M_{}`$ (which corresponds to $`21<M_B<18`$) is approximately in agreement with our result and P00, but their models cannot fully reproduce the drop in the merger fraction by a factor of 2 for the brighter galaxies. The simple ’chance hypothesis’, where galaxies fall into each other’s gravitational influence zone simply by random motions, appears to work somewhat better: galaxies with $`22<M_B<19`$ have an approximately 6 times lower chance to be in close proximity than galaxies with $`21<M_B<18`$, by integration of the luminosity function, but such objects are also $`3`$ times more clustered (Norberg et al., 2002). Therefore the pair fraction should be lower by about a factor of 2 which is in reasonable agreement with our findings. This is apparently in contrast with the conclusions of Xu et al. (2004) who appear to reject the chance hypothesis for dwarfs in the proximity of giants; the hypothesis may only apply to giant galaxies, which are more resilient than dwarfs. Clearly, a more consistent interpretation of these results must await a more thorough theoretical modelling of this phenomenon.
Since most companions are found in pairs, rather than in triplets or higher order $`N`$-tuples, our derived value of $`N_c`$ is comparable to the fraction of galaxies in close pairs. Following Yee and Ellingson (1995) and P00 we now assume that half of our dynamically close pairs are actually physically close pairs and will hence merge. We now integrate the merger rate to $`z=1`$ to derive the fraction of present day galaxies that have undergone a major merger since this time. The fraction of merger remnants is (P00):
$$f_{rem}=1\underset{k=1}{\overset{N}{}}\frac{1F_{mg}(z_k)}{10.5F_{mg}(z_k)}$$
(10)
where $`F_{mg}(z)`$ is the merger rate at redshift $`z`$, which is assumed to vary as $`F_{mg}(z)=F_{mg}(0)(1+z)^m`$, with $`m=3`$, $`z_k`$ corresponds to a lookback time of $`t=kT_{mg}`$ and $`T_{mg}`$ is the merger timescale. Following P00, we assume that this merger timescale is 0.5 Gyr for galaxies with $`21<M_B<18`$ and, scaling by the luminosity of the galaxies, 0.2 Gyr for galaxies with $`22<M_B<19`$. We therefore have 27 merger timescales to $`z1`$ for the brighter sample and 11 timescales for the fainter sample. This implies that $`22.7\%`$ of $`22<M_B<19`$ galaxies and $`19.2\%`$ of $`21<M_B<18`$ galaxies have undergone a major merger since $`z1`$, approximately the last half of the Hubble time. P00 derive a remnant fraction $`6.6\%`$ for their sample, but they assume a slightly different cosmology and $`m=0`$, while Lin et al. (2004) also obtain a remnant fraction of $`10`$%, but use $`m=0.51`$ and P02 derived an integrated merger fraction of 15% with $`m=2.3`$. Note, however, that these values are strongly dependent on the assumed merger timescales, index for the evolution of the merger rate and fraction of merging galaxies. For the sake of comparison, $`\mathrm{\Lambda }`$CDM models predict that about 50% of $`L>L^{}`$ galaxies have undergone a major merger since $`z1`$ (Murali et al., 2002), which is about a factor of 2 lower than our estimate.
### 4.2 The properties of galaxies in pairs
It is interesting to compare the properties of galaxies in pairs with respect to their parent sample. Since only $`20\%`$ of these objects have undergone a merger in the last half of the Hubble time, this should provide a sample of ‘undisturbed’ galaxies and offer some insight as to the effects of close interactions (and mergers) on galaxy morphology and star formation. Figure 2 already offers a hint of this: many galaxies show obvious signs of interactions, with tidal tails, rings and other signs of morphological disturbance, a connection exploited by the asymmetry parameter (Conselice et al., 2003).
Figure 3 shows the colour distribution in rest $`ur`$ for galaxies in both samples, compared with that of their parent distribution as well as the distribution of morphologies in the pairs and parent samples. Although the errors on the distributions for paired galaxies are large, because of small number statistics, there is some evidence that galaxies in pairs tend to be bluer than those in the parent sample, by approximately 0.1 mag in $`ur`$, suggesting that close interactions induce star formation episodes (Alonso et al., 2004; Nikolic et al., 2004; Barton et al., 2003; Lambas et al., 2003; Barton et al., 2000), although the average blueing for our sample is more consistent with relatively mild starbursts (Bergvall et al., 2003). On the other hand, not all of our kinematic pairs may be real merging systems and therefore, not all objects may show enhanced star formation.
Comparison of the morphologies (determined from visual examination of the images) shows that there are marginally more E/S0 galaxies and Sd/Irr galaxies among pairs, while there are fewer Sabc galaxies. The excess of early-type galaxies in pairs was noted by Junqueira et al. (1998), and this was attributed to an excess of lenticulars by Rampazzo and Sulentic (1992) and Rampazzo et al. (1995) who suggest that pairs containing early-type members may be the result of early merging of groups (which would have additional companions handy to fuel further merging). The excess of very late-type galaxies may be explained by interactions increasing the star formation rate and therefore leading to later classifications. These may be long-lived interacting systems rather than objects in an early phase of a merger (Junqueira et al., 1998).
Mergers and interactions therefore alter the morphologies of the member galaxies and induce star formation episodes. Scenarios in which local conditions (i.e. mergers and close encounters) are primarily responsible for morphological evolution and the decline of star formation in denser environments have been recently proposed following analysis of large local samples of field and cluster galaxies (Balogh et al., 2004; De Propris et al., 2004; Croton et al., 2005), although the change in the fraction of blue galaxies is much higher than the merger fraction we derive (suggesting a large contribution from lower luminosity galaxies). Wider and deeper surveys would be valuable, in allowing consideration of the importance of minor mergers.
RDP is supported by a PPARC fellowship. JL acknowledges an ESO fellowship. The Millennium Galaxy Catalogue consists of imaging data from the Isaac Newton Telescope and spectroscopic data from the Anglo Australian Telescope, the ANU 2.3m, the ESO New Technology Telescope, the Telescopio Nazionale Galileo, and the Gemini Telescope. The survey has been supported through grants from the Particle Physics and Astronomy Research Council (UK) and the Australian Research Council (AUS). The data and data products are publicly available from http://www.eso.org/$``$jliske/mgc/ or on request from JL or SPD. We would like to thank the referee, D. R. Patton, for a number of helpful comments and suggestions that made the paper clearer and better.
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# The Jülich hyperon-nucleon model revisited
## I Introduction
The role of strangeness in low and medium energy nuclear physics is currently of considerable interest, as it has the potential to deepen our understanding of the relevant strong interaction mechanisms in the non-perturbative regime of QCD. For example, the system of a strange baryon (hyperon $`Y`$) and a nucleon ($`N`$) is in principle an ideal testing ground for investigating the importance of SU(3)<sub>flavor</sub> symmetry in hadronic interactions. Existing meson exchange models of the $`YN`$ force usually assume SU(3) symmetry for the hadronic coupling constants, and in some cases Holz ; Reu even the SU(6) symmetry of the quark model. The symmetry requirements provide relations between couplings of mesons of a given multiplet to the baryon current, which greatly reduce the number of free model parameters. Specifically, coupling constants at the strange vertices are then connected to nucleon-nucleon-meson coupling constants, which in turn are constrained by the wealth of empirical information on $`NN`$ scattering. Essentially all these $`YN`$ interaction models can reproduce the existing $`YN`$ scattering data, so that at present the assumption of SU(3) symmetry for the coupling constants cannot be ruled out by experiment.
One should note, however, that the various models differ dramatically in the treatment of the scalar-isoscalar meson sector, which describes the baryon-baryon interaction at intermediate ranges. For example, the Nijmegen group NijII ; NijIII ; NijIV ; NijV views this interaction as being generated by genuine scalar meson exchange. In their model D NijII an $`ϵ(760)`$ is exchanged as an SU(3)<sub>flavor</sub> singlet. In models F NijIII , NSC NijIV , and NSC97 NijV a scalar SU(3) nonet is exchanged — namely, two isospin-0 mesons (besides the $`ϵ(760)`$, the $`ϵ^{}(1250)`$ in model F and $`S^{}(975)`$ ($`f_0(980)`$) in model NSC (NSC97)), an isospin-1 meson ($`\delta `$ or $`a_0(980)`$) and an isospin-1/2 strange meson $`\kappa `$ with a mass of 1000 MeV. A genuine scalar SU(3) nonet is also present in the so-called Ehime potential Ehime , where besides the $`S^{}(975)`$ and $`\delta `$ (or $`a_0(980)`$) the $`f_0(1581)`$ and the $`K_0^{}(1429)`$ are included. In additon the model incorporates two effective scalar-meson exchanges, $`\sigma (484)`$ and $`\kappa (839)`$, that stand for $`(\pi \pi )_{I=0}`$ and $`(K\pi )_{I=1/2}`$ correlations but are treated phenomenologically. The Tübingen model, on the other hand, which is essentially a constituent quark model supplemented by $`\pi `$ and $`\sigma `$ exchange at intermediate and short ranges, treats the $`\sigma `$ meson as an SU(3) singlet with a mass of 520 MeV Tueb or 675 MeV Zhang1 , respectively. Finally, in the quark-models of Zhang et al. Zhang2 and Fujiwara et al. Fujiwara a scalar SU(3) nonet is exchanged, though in this case between quarks and not between the baryons.
In the (full) Bonn $`NN`$ potential MHE the intermediate range attraction is provided by uncorrelated and correlated $`\pi \pi `$ exchange processes (Figs. 1(a)–(b) and Fig. 1(c), respectively), with $`NN`$, $`N\mathrm{\Delta }`$ and $`\mathrm{\Delta }\mathrm{\Delta }`$ intermediate states. From earlier studies of the $`\pi \pi `$ interaction it is known that $`\pi \pi `$ correlations are important mainly in the scalar-isoscalar and vector-isovector channels. In one-boson-exchange (OBE) potentials these are included effectively via exchange of sharp mass $`\sigma `$ and $`\rho `$ mesons. One disadvantage of such a simplified treatment is that this parameterization cannot be transferred into the hyperon sector in a well defined manner. Therefore in the earlier $`YN`$ interaction models of the Jülich group Holz , which start from the Bonn $`NN`$ potential, the coupling constants of the fictitious $`\sigma `$ meson at the strange vertices ($`\mathrm{\Lambda }\mathrm{\Lambda }\sigma `$, $`\mathrm{\Sigma }\mathrm{\Sigma }\sigma `$) are free parameters — a rather unsatisfactory feature of the models. This is especially true for the extension to the strangeness $`S=2`$ channels, interest in which initiated with the prediction of the H-dibaryon by Jaffe Jaffe .
These problems can be overcome by an explicit evaluation of correlated $`\pi \pi `$ exchange in the various baryon-baryon channels. A corresponding calculation was initially done only for the $`NN`$ case (Fig. 1(c)) in Ref. Kim , but was extended in a recent paper REUBER by the Jülich group so that now a full and consistent microscopic derivation of correlated $`\pi \pi `$ exchange in various baryon-baryon ($`BB^{}`$) channels with strangeness $`S=0,1`$ and $`2`$ is available. The starting point was a field theoretical model for both the $`N\overline{N}\pi \pi `$ Born amplitudes and the $`\pi \pi `$ and $`K\overline{K}`$ elastic scattering Lohse ; Janssen ; Schutz . Thus, the $`K\overline{K}`$ channel is treated on an equal footing with the $`\pi \pi `$ channel in order to reliably determine the influence of $`K\overline{K}`$ correlations in the relevant $`t`$-channels. Then, with the help of unitarity and dispersion relations the amplitude for the correlated $`\pi \pi `$ exchange in the $`NN`$ channel but also for the $`YN`$ and $`YY`$ systems were computed. Thus, within this approach one can replace the phenomenological $`\sigma `$ and $`\rho `$ exchanges in the Bonn $`NN`$ MHE and Jülich $`YN`$ Holz models by correlated processes, i.e. eliminate undetermined parameters such as the $`BB^{}\sigma `$ coupling constants.
In the present paper a new $`YN`$ model is presented that utilizes this microscopic model of correlated $`\pi \pi `$ and $`K\overline{K}`$ exchange to fix the contributions in the scalar-isoscalar ($`\sigma `$) and vector-isovector ($`\rho `$) channels. The model incorporates also the standard one boson exchange contributions of the lowest pseudoscalar and vector meson multiplets with coupling constants determined by SU(6) symmetry relations. Assuming the SU(6) symmetry means that also the so-called $`F/(F+D)`$ ratios are fixed. In addition, there are further new ingrediations as compared to the original Jülich $`YN`$ model Holz . First of all, the contribution from the $`a_0(980)`$ meson is taken into account. Secondly, we consider the exchange of a strange scalar meson, the $`\kappa `$, with mass $`1000`$ MeV. Let us emphasize, however, that in analogy with the $`\sigma `$ meson these particles are likewise not viewed as being members of a scalar meson SU(3) multiplet, but rather as representations of strong meson-meson correlations in the scalar–isospin-1/2 ($`\pi K`$) Lohse and scalar–isovector ($`\pi \eta `$$`K\overline{K}`$) Janssen channels respectively. In principle, their contributions can also be evaluated along the lines of Ref. REUBER , however, for simplicity in the present model they are effectively parameterized by one-boson-exchange diagrams with the appropriate quantum numbers assuming the coupling constants to be free parameters.
In the next two sections we describe the principal steps of the derivation of correlated $`\pi \pi `$ and $`K\overline{K}`$ exchange potentials for the baryon–baryon amplitudes in the $`\sigma `$ and $`\rho `$ channels. In particular, in Sect. 2 we give a short outline of the microscopic model for the required $`B\overline{B^{}}\pi \pi ,K\overline{K}`$ amplitudes. The derivation of the potentials themselves is indicated in Section 3. Furthermore, we introduce and discuss the parameterization of correlated $`\pi \pi `$ and $`K\overline{K}`$ exchange potentials by an effective $`\sigma `$ and $`\rho `$ exchange for the $`YN`$ channels. These effective parameterizations are then adopted for the construction of the new $`YN`$ model. In Sect. 4 the other ingredients of our $`YN`$ model are introduced. Specifically, we comment on the employed strategy for fixing the parameters of the model. Then we present and discuss numerical results of the model for $`YN`$ scattering observables, phase shifts and effective range parameters. Finally, some concluding remarks are made in Section 5.
## II Model for correlated $`2\pi `$ exchange
Based on a $`\pi \pi K\overline{K}`$ amplitude the evaluation of diagrams such as in Fig. 1(c) for any $`BB^{}`$ system can be done in two steps. Firstly the $`N\overline{N}(\mathrm{\Lambda }\overline{\mathrm{\Lambda }},\mathrm{\Sigma }\overline{\mathrm{\Sigma }},\mathrm{etc}.)2\pi ,K\overline{K}`$ amplitudes are determined in the pseudophysical region ($`t4m_\pi ^2`$) and then dispersion theory and unitarity are applied to connect those amplitudes with the corresponding physical amplitudes in the various baryon-baryon channels.
Figure 2 shows a graphic representation of our dynamical model for correlated $`2\pi K\overline{K}`$ exchange. Here $`B\overline{B^{}}`$ stands for $`N\overline{N}`$, $`\mathrm{\Lambda }\overline{\mathrm{\Lambda }}`$, $`\mathrm{\Lambda }\overline{\mathrm{\Sigma }}`$/$`\mathrm{\Sigma }\overline{\mathrm{\Lambda }}`$ or $`\mathrm{\Sigma }\overline{\mathrm{\Sigma }}`$. Formally the amplitudes for the processes $`B\overline{B^{}}\alpha `$ (with $`\alpha =\pi \pi `$, $`K\overline{K}`$) are obtained from solving the scattering equation
$$T_{B,\overline{B^{}}\alpha }=V_{B,\overline{B^{}}\alpha }+\underset{\beta =\pi \pi ,K\overline{K}}{}T_{\alpha ,\beta }G_\beta V_{B,\overline{B^{}}\beta }.$$
(1)
Here $`T_{\alpha ,\beta }`$ is the $`\pi \pi K\overline{K}`$ (coupled-channel) reaction amplitude, $`V_{B,\overline{B^{}}\beta }`$ the $`B\overline{B^{}}\pi \pi ,K\overline{K}`$ transition Born amplitude and $`G_\beta `$ the free ($`\pi \pi `$ or $`K\overline{K}`$) Green’s function The first two quantities are the basic ingredients of the model. These amplitudes have to be known in the so-called pseudophysical region, i.e. for energies below the $`B\overline{B^{}}`$ threshold. While for $`N\overline{N}\pi \pi `$ the corresponding amplitudes can be derived from empirical information on $`\pi N`$ and $`\pi \pi `$ scattering via an analytic continuation, this is not possible for the transitions $`Y\overline{Y^{}}\pi \pi ,K\overline{K}`$. Thus, a microscopic model for the $`B\overline{B^{}}\pi \pi `$,$`K\overline{K}`$ is a pre-requisite for the evaluation of the correlated $`\pi \pi `$ and $`K\overline{K}`$ exchange in the $`YN`$ and $`YY`$ channels. Such a model was constructed by Reuber et al. in Ref. REUBER . A main feature of this model is the completely consistent treatment of its two components, namely the $`B\overline{B^{}}\pi \pi ,K\overline{K}`$ Born amplitudes and the $`\pi \pi `$-$`K\overline{K}`$ correlations with respect to the $`\sigma `$\- and $`\rho `$-channels. Both components are derived in field theory from an ansatz for the hadronic Lagrangians REUBER . The considered contributions are briefly described in the subsections below. The subsequent evaluation of the baryon-baryon interaction via dispersion theory and unitarity is summarized in the next section.
### II.1 The $`\pi \pi K\overline{K}`$ amplitudes
The dynamical model used for the $`\pi \pi K\overline{K}`$ amplitudes is derived within the meson exchange framework and involves the $`\pi \pi `$ and $`K\overline{K}`$ coupled channels Lohse ; Janssen ; Schutz . The driving terms for the diagonal interactions consist of ($`t`$-channel) exchange diagrams ($`\rho `$ and $`\rho `$, $`\omega `$, $`\varphi `$, respectively) and ($`s`$-channel) pole diagrams with $`ϵf_0(1440)`$, $`\rho \rho (770)`$ and $`f_2f_2(1274)`$ intermediate states. The coupling $`\pi \pi K\overline{K}`$ is provided by $`K^{}(892)`$ exchange. The corresponding diagrams are shown in Fig. 3. The potentials derived from those diagrams are iterated in a coupled-channel Lippmann-Schwinger-type scattering equation. The free parameters of the $`\pi \pi K\overline{K}`$ model were adjusted to the empirical $`\pi \pi `$ phase shifts and inelasticities. For details on the model and a comparision of the resulting $`\pi \pi `$ phase shifts with experimental values we refer the reader to Refs. Janssen ; Schutz .
### II.2 The $`B\overline{B}2\pi ,K\overline{K}`$ Born amplitudes
The Born amplitudes for the transition $`B\overline{B}\alpha `$ with $`\alpha =\pi \pi ,K\overline{K}`$ are built up from an ($`s`$-channel) $`\rho `$-pole diagram and all possible diagrams involving the exchange of baryons out of the $`J^P=\frac{1}{2}^+`$ octet or the $`J^P=\frac{3}{2}^+`$ decuplet REUBER . For illustration we show in Fig. 4 those diagrams that contribute to the transition amplitude for $`\mathrm{\Sigma }\overline{\mathrm{\Sigma }}2\pi ,K\overline{K}`$.
In the construction of the model the number of free parameters has been kept to a minimum. Specifically, the coupling constants at the various vertices involving the pseudoscalar mesons were fixed by SU(6) symmetry relations. As far as the $`\rho `$-pole diagram is concerned the (bare) coupling constants and form factors at the $`\pi \pi \rho ^{(0)}`$ and $`K\overline{K}\rho ^{(0)}`$ vertices were already determined in the model for the $`\pi \pi K\overline{K}`$ interaction Schutz and were taken over from there. Then, assuming that the bare $`\rho `$-meson couples universally to the isospin current all vector couplings $`g_{BB^{}\rho }^{(0)}`$ to the baryonic vertices were fixed as well. For the tensor couplings $`f_{BB^{}\rho }^{(0)}`$ again SU(6) symmetry relations were applied.
The four remaining free parameters (the tensor coupling constant $`f_{NN\rho }^{(0)}`$, the parameter $`x_\mathrm{\Delta }`$ characterizing the strength of the off-shell part in the $`\mathrm{\Delta }N\pi `$ Lagrangian<sup>1</sup><sup>1</sup>1In a more modern language, it can be shown that such off-shell parameters really correspond to low-energy constants of non-propagating contact interactions, see e.g. BKMdelta , and form-factor parameters for the exchanged baryons for the octet and decuplet, respectively REUBER ) were fixed by adjusting the model predictions to the quasi-empirical information on the amplitudes $`N\overline{N}\pi \pi `$ obtained by Höhler at el. Hoehler2 by analytically continuing the $`\pi N`$ and $`\pi \pi `$ scattering data. Once this is done the model can be used to generate the amplitudes for any $`B\overline{B}^{}\pi \pi ,K\overline{K}`$ channel. Though it should have become clear already from the discussion above, we want to emphasize again that the extrapolation of the model for the $`N\overline{N}\pi \pi `$ amplitudes to other channels depends crucially on the assumption of SU(3) symmetry for the pseudo-scalar sector and it is based also on the hope that the correct description of the quasiempirical $`N\overline{N}\pi \pi `$ amplitudes guarantees a reasonable description of the other baryon-antibaryon channels, for which no empirical information is available.
## III Potential from correlated $`\pi \pi `$ and $`K\overline{K}`$ exchange
Assuming analyticity for the amplitudes dispersion relations can be formulated for the baryon-baryon amplitudes, which connect physical amplitudes in the $`s`$-channel with singularities and discontinuities of these amplitudes in the pseudophysical region of the $`t`$-channel processes for the $`J^P=0^+`$ ($`\sigma `$) and $`1^{}`$ ($`\rho `$) channel:
$$V_{B_1,B_2B_1^{},B_2^{}}^{(0^+,1^{})}(t)_{4m_\pi ^2}^{\mathrm{}}𝑑t^{}\frac{\mathrm{Im}V_{B_1,\overline{B_1^{}}\overline{B_2},B_2^{}}^{(0^+,1^{})}(t^{})}{t^{}t},t<0.$$
(2)
Via unitarity relations the singularity structure of the baryon-baryon amplitudes for $`\pi \pi `$ and $`K\overline{K}`$ exchange are fixed by and can be written as products of the $`B\overline{B^{}}\pi \pi ,K\overline{K}`$ amplitudes
$$\mathrm{Im}V_{B_1,\overline{B_1^{}}\overline{B_2},B_2^{}}^{(0^+,1^{})}(t^{})\underset{\alpha =\pi \pi ,K\overline{K}}{}T_{B_1,\overline{B_1^{}}\alpha }^{,(0^+,1^{})}T_{\overline{B_2},B_2^{}\alpha }^{(0^+,1^{})}.$$
(3)
Thus, from the $`B\overline{B^{}}2\pi `$ helicity amplitudes the spectral functions can be calculated
$$\rho _{B_1,B_2B_1^{},B_2^{}}^{(0^+,1^{})}(t^{})\underset{\alpha =\pi \pi ,K\overline{K}}{}T_{B_1,\overline{B_1^{}}\alpha }^{,(0^+,1^{})}T_{\overline{B_2},B_2^{}\alpha }^{(0^+,1^{})}$$
(4)
which are then inserted into dispersion integrals to obtain the (on-shell) baryon-baryon interaction:
$$V_{B_1,B_2B_1^{},B_2^{}}^{(0^+,1^{})}(t)_{4m_\pi ^2}^{\mathrm{}}𝑑t^{}\frac{\rho _{B_1,B_2B_1^{},B_2^{}}^{(0^+,1^{})}(t^{})}{t^{}t},t<0.$$
(5)
The underlying formalism is quite involved and has been outlined in detail already in Ref. REUBER . Thus, we refrain from repeating it here. Rather we want to provide only some more general information.
Since the dispersion-theoretical evaluation is restricted to the contribution of (correlated) $`\pi \pi `$ and $`K\overline{K}`$ exchange to the baryon-baryon amplitudes only those singularities are taken into account which are generated by $`\pi \pi `$ and $`K\overline{K}`$ intermediate states, namely the discontinuities due to the $`\pi \pi `$ and $`K\overline{K}`$ unitarity cut (the so-called right-hand cut). The left-hand cuts, which are due to unitarity constraints for the $`u`$-channel reaction, can be neglected in the baryon-baryon channels considered here, since they start at large, negative $`t`$-values (from which they extend to $`\mathrm{}`$) and are therefore far away from the physical region relevant for low-energy $`s`$-channel processes.
The $`B\overline{B^{}}\alpha `$ amplitudes, which enter in Eq. (3) are derived from a microscopic model which is based on the hadron-exchange picture, cf. Sect. II. Of course, this model has a limited range of validity: for energies far beyond $`t_{max}^{}100m_\pi ^2`$ it cannot provide reliable results. The dispersion integral for the invariant amplitudes extending in principle along the whole $`\pi \pi `$ right-hand cut has therefore to be limited to an upper bound, $`t_{max}^{}`$, which has been put to $`t_{max}^{}`$ = 120 $`m_\pi ^2`$ in Ref. REUBER .
The spectral function (4) for the ($`0^+`$) $`\sigma `$-channel has only one component but the one for the ($`1^{}`$) $`\rho `$-channel consists of four linearly independent components, which reflects the more complicated spin structure of this channel.
Finally, we should note that the helicity amplitudes obtained according to Fig. 2 still generate the uncorrelated (first diagram on the r.h.s. of Fig. 2), as well as the correlated pieces (second and third diagrams). Thus, in order to obtain the contribution of the truely correlated $`\pi \pi `$ and $`K\overline{K}`$ exchange one must eliminate the former from the spectral function. This is done by calculating the spectral function generated by the Born term and subtracting it from the total spectral function:
$$\rho ^{(0^+,1^{})}\rho ^{(0^+,1^{})}\rho _{\mathrm{Born}}^{(0^+,1^{})}.$$
(6)
In practice this means that, e.g., for the full Bonn $`NN`$ model contributions involving spin-1/2 as well as spin-3/2 baryons have to be subtracted since corresponding contributions are already treated explicitly in the s-channel in this model, namely via box diagrams with intermediate $`\mathrm{\Delta }`$-states as shown in Fig. 1(a). On the other hand only uncorrelated contributions involving spin-1/2 baryons are to be subtracted from the discontinuities of the invariant baryon-baryon amplitudes in order to avoid double counting if a simple OBE-model is used in the $`s`$-channel. This is the relevant procedure for the $`YN`$ model that will be presented in the next section.
Note that the spectral functions characterize both the strength and range of the interaction. Clearly, for sharp mass exchanges the spectral function becomes a $`\delta `$-function at the appropriate mass.
For convenience the authors of Ref. REUBER have presented their results in terms of effective coupling strengths, by parameterizing the correlated processes by (sharp mass) $`\sigma `$ and $`\rho `$ exchanges. The interaction potential resulting from the exchange of a $`\sigma `$ meson with mass $`m_\sigma `$ between two $`J^P=1/2^+`$ baryons $`A`$ and $`B`$ has the structure:
$$V_{A,BA,B}^\sigma (t)=g_{AA\sigma }g_{BB\sigma }\frac{F_\sigma ^2(t)}{tm_\sigma ^2},$$
(7)
where a form factor $`F_\sigma (t)`$ is applied at each vertex, taking into account the fact that the exchanged $`\sigma `$ meson is not on its mass shell. This form factor is parameterized in the conventional monopole form,
$$F_\sigma (t)=\frac{\mathrm{\Lambda }_\sigma ^2m_\sigma ^2}{\mathrm{\Lambda }_\sigma ^2t},$$
(8)
with a cutoff mass $`\mathrm{\Lambda }_\sigma `$ assumed to be the same for both vertices. The correlated potential as given in Eq. (2) can now be parameterized in terms of $`t`$-dependent strength functions $`G_{B_1^{},B_2^{}B_1,B_2}(t)`$, so that for the $`\sigma `$ case:
$$V_{A,BA,B}^{(0^+)}(t)=G_{ABAB}^\sigma (t)F_\sigma ^2(t)\frac{1}{tm_\sigma ^2}.$$
(9)
The effective coupling constants are then defined as:
$$g_{AA\sigma }g_{BB\sigma }G_{ABAB}^\sigma (t)=\frac{(tm_\sigma ^2)}{\pi F_\sigma ^2(t)}_{4m_\pi ^2}^{\mathrm{}}\frac{\rho _{ABAB}^{(0^+)}(t^{})}{t^{}t}𝑑t^{}.$$
(10)
Similar relations can be also derived for the correlated exchange in the isovector-vector channel REUBER , which in this case will involve vector as well as tensor coupling pieces.
It should be stressed that, so far, this parameterization does not involve any approximations as long as the full $`t`$-dependence of the effective coupling strengths is taken into account. The parameters of the $`\sigma `$ and $`\rho `$ exchange have been chosen to have the same values in all particle channels. The masses $`m_\sigma `$ and $`m_\rho `$ of the exchanged particles have been set to the values used in the Bonn-Jülich models of the $`NN`$ MHE and $`YN`$ Holz interactions, $`m_\sigma =550`$ MeV, $`m_\rho =770`$ MeV. The cutoff masses $`\mathrm{\Lambda }_\sigma `$ and $`\mathrm{\Lambda }_\rho `$ have been chosen so that the coupling strengths in the $`S=0,1`$ baryon-baryon channels vary only weakly with $`t`$. The resulting values ($`\mathrm{\Lambda }_\sigma =2.8`$ GeV, $`\mathrm{\Lambda }_\rho =2.5`$ GeV) are quite large compared to the values of the phenomenological parameterizations used in Refs. Holz ; MHE , and thus represent rather hard form factors.
Note that in the OBE framework the contribution of a genuine (SU(3)) $`\sigma `$ meson to the three reactions $`NNNN`$, $`YNYN`$, $`YYYY`$ is determined by two parameters (coupling constants), namely $`g_{NN\sigma }`$ and $`g_{YY\sigma }`$, whereas the correlated exchange is characterized by three independent strength functions ($`G_{NNNN}`$, $`G_{YNYN}`$, $`G_{YYYY}`$) so that vertex coupling constants cannot be determined uniquely. This implies directly that the strength parameters cannot fulfill SU(3) relations.
In the physical region the strength of the contributions is to a large extent governed by the value of $`G`$ at $`t=0`$. Those values for the various channels were tabulated in Ref. REUBER (cf. Tables 5-7) for the case of the full model calculation and also when uncorrelated contributions involving spin-1/2 baryons only are subtracted from the spectral function of the invariant baryon-baryon amplitudes. The latter are the proper values to be used for constructing a $`YN`$ model based on simple OBE-exchange diagrams.
In principle, the average size of the effective coupling strengths is only an approximate measure of the strength of correlated $`\pi \pi `$ and $`K\overline{K}`$ exchange in the various particle channels. The precise energy dependence of the correlated exchange as well as its relative strength in the different partial waves of the $`s`$-channel reaction is determined by the spectrum of exchanged invariant masses, or spectral functions, leading to a different $`t`$-dependence of the effective coupling strengths. This was demonstrated in Ref. Melni1 where the on-shell $`NN`$, $`\mathrm{\Lambda }N`$ and $`\mathrm{\Sigma }N`$ potentials in spin-singlet states with angular momentum $`L=0,2`$ and 4, generated directly by the scalar-isoscalar part of the correlated $`\pi \pi `$ and $`K\overline{K}`$ exchange, were compared to the corresponding results based on a $`\sigma `$ exchange with sharp mass. It could be seen that the correlated $`2\pi `$ exchange is significantly stronger in high partial waves because the $`\sigma `$ exchange, which corresponds to a spectral function proportional to $`\delta (t^{}m_\sigma ^2)`$, does not contain the long-range part of the correlated processes. Thus, parameterizing the results derived from the microscopic model by $`\sigma `$ exchange with a sharp mass, but using the effective coupling strength $`G_{NNNN}^\sigma `$ at $`t=0`$ one can obtain rough agreement with the exact result in the $`S`$ waves, say, but usually underestimates the magnitude considerably in the high partial waves. Obviously the replacement of correlated $`\pi \pi `$ and $`K\overline{K}`$ exchanges by an exchange of a sharp mass $`\sigma `$ meson with a $`t`$-independent coupling cannot provide a simultaneous description of both low and high partial waves.
These features are important for investigations of the $`NN`$ systems where the phase shifts are known quantitatively even for rather high partial waves. In this case the results of the correlated exchange should be used directly Kim . However, for the $`\mathrm{\Lambda }N`$ and $`\mathrm{\Sigma }N`$ systems only scattering observables are available, and those (total and differential cross sections) are primarily sensitive to $`S`$\- and $`P`$-wave contributions. Thus, here it is reasonable to simplify the calculation and use only an effective parametrization of the results derived from the microscopic model in terms of a $`\sigma `$ and $`\rho `$ exchange with a sharp mass. Specifically, combining Eqs. (7) and (9) we use the expression
$$V_{A,BA,B}^{(0^+)}(t)=G_{ABAB}^\sigma \stackrel{~}{F}_\sigma ^2(t)\frac{1}{tm_\sigma ^2}.$$
(11)
with
$$\stackrel{~}{F}_\sigma (t)=\frac{\mathrm{\Lambda }_\sigma ^2}{\mathrm{\Lambda }_\sigma ^2t}$$
(12)
and a similar one for the $`\rho `$ exchange contribution. The effective coupling strength $`G_{YNY^{}N}^\sigma `$ (and $`{}_{}{}^{ij}G_{YNY^{}N}^{\rho }`$) is deduced via Eq. (10) (and via a similar one for the $`\rho `$ channel, cf. Ref. REUBER ) for the form factor (12) and adjusted to the value at $`t=0`$. The different prescription for the vertex form factor as compared to Ref. REUBER , i.e to Eq. (8), is adopted here because it guarantees that the on-shell behaviour of the potential (which is fully determined by the dispersion integral) is not modified strongly as long as the energy is not too high. At the same time smaller cutoff masses as those mentioned above (and employed in REUBER ) can be used to ensure sufficient convergence when the potential (11) is iterated in the scattering equation. The concrete values used for the cutoff masses are $`\mathrm{\Lambda }_\sigma `$ = 2.5 (1.6) GeV for the $`\mathrm{\Lambda }N`$ ($`\mathrm{\Sigma }N`$) channels and $`\mathrm{\Lambda }_\rho `$ = 1.25 (1.8) GeV for the $`\mathrm{\Lambda }N\mathrm{\Sigma }N`$ transition ($`\mathrm{\Sigma }N`$ channel).
The effective coupling strengths employed in our new $`YN`$ model are compiled in Table 1. Though these values differ slightly from those given in Tables 5-7 of Ref. REUBER , due to the different choice of the form factor, we would like to emphasize that the strengths of the interactions at $`t=0`$ are the same in both cases and coincide with the one derived from the microscopic model of $`\pi \pi `$ and $`K\overline{K}`$ correlations.
In order to demonstrate that, we show in Fig. 5 the corresponding on-shell potential matrix elements for the $`{}_{}{}^{1}S_{0}^{}`$ partial wave of the $`\mathrm{\Lambda }N`$ and $`\mathrm{\Sigma }N`$ channels. One can see that in case of the $`\mathrm{\Lambda }N`$ system the result generated by the scalar-isoscalar part of correlated $`\pi \pi `$ and $`K\overline{K}`$ exchange is similar to the one of the $`\sigma `$ exchange used in the Jülich $`YN`$ model A. In fact, correlated $`\pi \pi `$ exchange is marginally stronger. It is also obvious that the parameterization of the interaction generated by correlated $`\pi \pi `$ exchange by an effective $`\sigma `$ exchange, c.f. the dotted line, works rather well. From the corresponding results for the on-shell $`\mathrm{\Sigma }N`$ potential one can see that here the $`\sigma `$ exchange used in the Jülich $`YN`$ model A is clearly much stronger than what one obtains from the correlated $`\pi \pi `$ and $`K\overline{K}`$ exchange. Once again the parameterization by an effective $`\sigma `$ exchange provides an excellent representation of the interaction strength.
## IV Results and discussion
### IV.1 Coupling constants
In the present $`YN`$ model we take into account exchange diagrams involving the well-established lowest lying pseudoscalar and vector meson SU(3) octets. Following the philosophy of the original Jülich $`YN`$ potential Holz the coupling constants in the pseudoscalar sector are fixed by strict SU(6) symmetry. In any case, this is also required for being consistent with the model of correlated $`\pi \pi `$ and $`K\overline{K}`$ exchange. The cutoff masses of the form factors belonging to the $`NN`$ vertices are taken over from the full Bonn $`NN`$ potential. The cutoff masses at the strange vertices are considered as open parameters though, in practice, their values are kept as close as possible to those found for the $`NN`$ vertices, cf. Table 2. Note that like in Holz and in line with the arguments brought forth in Ref. Reu we neglect again the contribution from $`\eta `$ meson exchange. Anyhow, in the full Bonn $`NN`$ model the $`\eta NN`$ coupling constant was set to zero. In addition phenomenological analyses Grein and also microscopic calculations, like those based on the topological chiral soliton model (extended Skyrme model with vector mesons) etaNN , indicate that this coupling constant should be small. Thus, the $`\eta `$ contribution would be completely unimportant anyway, given the pseudoscalar nature of its coupling. For the same reason the $`\eta ^{}`$ contribution is likewise not considered.
In the vector meson sector we depart from the strategy of the original Jülich $`YN`$ potential. As already mentioned above, first and most importantly the contribution of the $`\rho `$ meson is no longer seen as resulting from the exchange of a genuine particle that belongs to the SU(3) vector meson octet but is identified with the strength generated by a microscopic model of correlated $`\pi \pi `$ and $`K\overline{K}`$ in the vector-isovector channel. The effective coupling constants for $`\rho `$ exchange in the various $`YN`$ and $`YY`$ channels have been extracted and thoroughly analysed in Ref. REUBER . Thereby it was found that the result from correlated exchange deviates significantly from those implied by SU(3) symmetry – even though SU(3) symmetry was imposed for the bare $`\rho NN`$ and $`\rho YY`$ couplings, cf. REUBER . In view of this it is questionable whether one should invoke SU(3) symmetry for fixing the other coupling strengths of the vector-meson octet, i.e. those of the $`K^{}`$ meson and of the coupling of the $`\omega `$ meson to the hyperons. But in absence of any better alternative we still follow this prescription for the present model. As reference values we take here the $`NN\rho `$ coupling constants of the full Bonn $`NN`$ potential MHE , which were already used for the old $`YN`$ model Holz ; Reu . However, as far as the $`YY\omega `$ coupling constants are concerned now we take into account the insight gained in Ref. Janssen1 that the $`\omega `$ exchange in the full Bonn $`NN`$ potential represents not only the genuine SU(3) $`\omega `$ but is also an effective parametrization of additional short-range contributions from correlated $`\pi \rho `$ exchange, say, that are not included explicitly in that model. Therefore, in the Bonn $`NN`$ model the required $`NN\omega `$ coupling constant is indeed much larger than what follows from the SU(3) relations and this large coupling constant formed also the basis for fixing the $`YY\omega `$ coupling constants of the old Jülich $`YN`$ model Holz ; Reu , cf. the discussion in Sect. 2.2 of Ref. Reu . In the present model we adopt the smaller value found in Ref. Janssen1 which is very close to the SU(3) value. This is in line with results obtained from a dispersion-theoretical analysis of the nucleon electromagnetic form factors - the inclusion of the $`\pi \rho `$ continuum sizeably reduces the $`\omega NN`$ coupling, compare the values found in MMD with the ones in MMSvO . Assuming furthermore that the $`\rho `$ meson couples universally to the isospin current – which fixes the $`F/(F+D)`$ ratio $`\alpha _V^e`$ to 1 – and ideal mixing for the $`\varphi `$ and $`\omega `$ mesons then yields the following relation for the $`\omega `$ coupling constants:
$`g_{\mathrm{\Lambda }\mathrm{\Lambda }\omega }=g_{\mathrm{\Sigma }\mathrm{\Sigma }\omega }={\displaystyle \frac{2}{3}}g_{NN\omega },f_{\mathrm{\Lambda }\mathrm{\Lambda }\omega }={\displaystyle \frac{5}{6}}f_{NN\omega }{\displaystyle \frac{1}{2}}f_{NN\rho },f_{\mathrm{\Sigma }\mathrm{\Sigma }\omega }={\displaystyle \frac{1}{2}}f_{NN\omega }+{\displaystyle \frac{1}{2}}f_{NN\rho }`$ (13)
For $`f_{NN\omega }`$ and f<sub>NNρ</sub> we take over the values of the full Bonn $`NN`$ potential. Since $`f_{NN\omega }`$=0 MHE it follows that $`f_{\mathrm{\Lambda }\mathrm{\Lambda }\omega }=f_{\mathrm{\Sigma }\mathrm{\Sigma }\omega }`$.
The short-range contributions from correlated $`\pi \rho `$ exchange were parametrized by an effective $`\omega ^{}`$ exchange in Ref. Janssen1 with a mass of $`m_\omega ^{}`$ = 1120 MeV. We follow here the same strategy but treat the coupling constants of the $`\omega ^{}`$ to the strange baryons as free parameters to be determined in a fit to the $`YN`$ data.
Like the $`\rho `$ also the contribution of the $`\sigma `$ meson is computed from a microscopic model of correlated $`\pi \pi `$ and $`K\overline{K}`$ – now from the scalar-isoscalar channel. The effective coupling constants for $`\sigma `$ exchange in the various $`YN`$ channels have been discussed in the previous section.
Besides replacing the conventional $`\sigma `$ and $`\rho `$ exchanges by correlated $`\pi \pi `$ and $`K\overline{K}`$ exchange, there are in addition some other new ingredients in the present $`YN`$ model. First of all, we now take into account contributions from $`a_0(980)`$ exchange. The $`a_0`$ meson is present in the original Bonn $`NN`$ potential MHE , and for consistency should also be included in the $`YN`$ model. Secondly, we consider the exchange of a strange scalar meson, the $`\kappa `$, with mass $`1000`$ MeV. Let us emphasize, however, that like in case of the $`\sigma `$ meson these particles are not viewed as being members of a scalar meson SU(3) multiplet, but rather as representations of strong meson-meson correlations in the scalar–isovector ($`\pi \eta `$$`K\overline{K}`$) Janssen and scalar–isospin-1/2 ($`\pi K`$) channels Lohse , respectively. In principle, their contributions can also be evaluated along the lines of Ref. REUBER , however, for simplicity in the present model they are effectively parameterized by one-boson-exchange diagrams with the appropriate quantum numbers assuming the coupling constants to be free parameters. The parameters specifying those ingredients are summarized in Table 3.
Thus we have the following scenario: The long- and intermediate-range part of our new $`YN`$ interaction model is completely determined by SU(6) constraints (for the pseudoscalar and to some extent also for the vector mesons) and by correlated $`\pi \pi `$ and $`K\overline{K}`$ exchange. The short-range part is viewed as being also due to correlated meson-meson exchanges but in practice is parametrized phenomelogically in terms of one-boson-exchange contributions in specific spin-isospin channels. In particular, no SU(3) relations are imposed on the short-range part. This assumption is based on our observation that the contributions in the $`\rho `$ exchange channel as they result from correlated $`\pi \pi `$ and $`K\overline{K}`$ no longer fulfill SU(3) relations, but it also acknowledges the fact that at present there is no general agreement about who are the actual members of the lowest-lying scalar meson SU(3) multiplet. A graphical representation of all meson-exchange contributions that are included in the new $`YN`$ model is given in Fig. 6.
In recent investigations of the $`NN`$ interaction within the framework of chiral perturbation theory Epelbaum only pionic degrees of freedom are taken into account and all short-range physics is parametrized by contact terms. This is certainly also an option that one should explore for the $`YN`$ system Korpa in the future Henk . As a first step we consider here an alternative model where the contributions of the $`\kappa `$(1000) meson - whose mass and even existence is still under dispute Kappa – are substituted by a contact term. In practice this means that we replace the product of the $`\kappa `$ coupling constants and propagator, $`G_{BB^{}BB^{}}/(m_\kappa ^2t)`$, by $`G_{BB^{}BB^{}}/m_\kappa ^2`$ and readjust only the parameters of related to the $`\kappa `$ exchange (with one exception). Those parameters can be found also in Table 3, in square brackets, for those cases where they differ from the values of our regular model. Results corresponding to the model with the contact term will also be presented in the next section.
In the fitting procedure we only take into account data on total cross sections (and energies near the corresponding thresholds) for the channels $`\mathrm{\Lambda }p`$ Alex ; Sechi ; Kadyk , $`\mathrm{\Sigma }^{}p`$ Eisele , $`\mathrm{\Sigma }^{}p\mathrm{\Lambda }n`$ Engel , $`\mathrm{\Sigma }^{}p\mathrm{\Sigma }^0n`$ Engel , and $`\mathrm{\Sigma }^+p`$ Eisele . Differential cross sections but also total cross sections at higher energies Stephen ; Kondo ; Ahn are therefore genuine predictions of our model. As already mentioned above, the free parameters in our model consist of the cut-off masses at the strange vertices and the coupling constants of the $`a_0`$(980), $`\kappa `$(1000) and $`\omega ^{}`$(1120) mesons. When adjusting those parameters to the empirical data it turned out that the results are not very sensitive to the cut-off masses in the pseudo-scalar sector and we fixed them to be close to the cutoff mass used at the $`\pi NN`$ vertex. There is also only a weak sensitivity to the cut-off masses used for the correlated $`\pi \pi `$-$`K\overline{K}`$ contributions in the $`\sigma `$ and $`\rho `$ channels. This is due to the chosen analytic form of the form factors that practically does not change the strength of the corresponding potentials as they result from the microscopic model – which is of course intended, cf. the discussion in Sect. III. Besides the cut-off masses of the vector mesons we found that also the parameters of the $`a_0`$(980) and $`\kappa `$(1000) mesons, viewed here as effective parametrization of correlated $`\pi \eta `$ and $`\pi K`$ exchange, have a sizeable influence. In fact, without the contributions of the latter two mesons we would not have been able to achieve a satisfactory description of the data. Note that those two exchanges were not considered in the original Jülich model Holz . We should say that values of the coupling strengths and cut-off masses for those scalar mesons are strongly correlated and cannot be fixed independently from a fit to the data. Thus, one should not attribute any physical significance to the actual values of the coupling strengths or cut-off masses that we found individually.
Finally we want to mention that the fit to the available $`YN`$ data did not constrain the relative magnitude of the $`{}_{}{}^{1}S_{0}^{}`$ and $`{}_{}{}^{3}S_{1}^{}`$ partial waves in the $`\mathrm{\Lambda }N`$ system. Thus, as a further constraint, we required the $`{}_{}{}^{1}S_{0}^{}`$ scattering length to be larger than the one for $`{}_{}{}^{3}S_{1}^{}`$ – as it seems to be necessary if one wants to achieve a bound hypertriton Ueda . A first application of the new $`YN`$ model in three-body calculations confirmed that it yields indeed a bound hypertriton state Nogga .
### IV.2 The scattering equation
The original $`YN`$ model of the Jülich group was derived within the framework of time ordered perturbation theory (TOPT) Holz . In this approach retardation effects from the meson-exchange diagrams are retained (and those of baryon-exchange as well) and as a consequence the interaction depends explicitly on the starting energy. This is not convenient if one wants to apply the $`YN`$ model in conventional few-body Miya1 ; Miya2 ; Miya4 ; Akaishi ; Nogga1 ; Nemura ; Fujiwara3N ; Hiyama or many-body Ramos ; Tzeng ; Fujii ; Lenske investigations. Thus, in Ref. Reu the Jülich group presented energy-independent versions of their $`YN`$ model where the energy dependence was removed in such a way that basically all other characteristics of the original model could be kept. The detailed comparison of the TOPT model and its energy-independent counterpart performed in Ref. Reu made clear that this goal was indeed achieved.
Since we are also interested to facilitate an application of our new $`YN`$ model in future few- and many-body investigations we will likewise present here an energy-independent interaction. This implies that we do not use the (relativistic) TOPT scattering equation of Ref. Holz but instead solve the nonrelativistic (coupled-channel) Lippmann-Schwinger equation
$$T_{i,j}=V_{i,j}+\underset{k}{}V_{i,k}G_kT_{k,j}$$
(14)
to obtain the scattering amplitude $`T_{i,j}`$. Here the indices ($`i,j,k`$) stand for the $`\mathrm{\Lambda }N`$ and $`\mathrm{\Sigma }N`$ channels and the nonrelativistic Green’s function $`G_k`$ is given by
$$G_k=\left[\frac{q_k^2𝐪^2}{2\mu _k}+i\epsilon \right]^1,$$
(15)
where $`\mu _k=M_YM_N/(M_Y+M_N)`$ is the reduced mass and $`𝐪^{}`$ the c.m. momentum in the intermediate $`YN`$ channel. $`q_k=q_k(z)`$ denotes the on-shell momentum in the intermediate $`YN`$ state defined by $`z=\sqrt{M_Y^2+q_k^2}+\sqrt{M_N^2+q_k^2}`$. The latter equation guarantees that the $`\mathrm{\Sigma }N`$ channel opens exactly at the physical threshold. Note that $`q_{\mathrm{\Sigma }N}`$ is imaginary for starting energies below the $`\mathrm{\Sigma }N`$ threshold ($`z<M_\mathrm{\Sigma }+M_N`$). Explicit expressions for the potential matrix elements $`V_{i,j}`$ for the various exchange diagrams can be found in Ref. Holz . The dependence on the starting energy $`z`$ is removed via the prescriptions given in Eq. (4.7) of Ref. Reu . Note that the potential matrix elements $`V_{i,j}`$ are derived by assuming isospin symmetry. However, the Lippmann-Schwinger equation (14) is solved in particle space using the proper physical masses of the baryons for the various $`\mathrm{\Sigma }N`$ channels. Furthermore, in the charged channels the Coulomb potential is taken into account. Since we solve the Lippmann-Schwinger equation in momentum space this is done by means of the Vincent-Phatak method Holz ; Phatak .
### IV.3 Hyperon-nucleon observables
In Fig. 7 we compare the integrated cross sections obtained from the new $`YN`$ potential (solid curves) with the $`YNY^{}N`$ scattering data. Obviously, a good reproduction of the empirical data Alex ; Sechi ; Kadyk ; Eisele ; Engel is achieved. Also shown are results from the original Jülich $`YN`$ model A Holz (dash-dotted curves). The main qualitative differences between the two models appear in the $`\mathrm{\Lambda }p\mathrm{\Lambda }p`$ channel, for which the Jülich model Holz (with standard $`\sigma `$ and $`\rho `$ exchange) predicts a broad shoulder at $`p_{lab}`$ 350 MeV/c. This structure, which is not supported by the available experimental evidence, is due to a bound state in the $`{}_{}{}^{1}S_{0}^{}`$ partial wave of the $`\mathrm{\Sigma }N`$ channel. It is not present in the new model anymore. (We should say, however, that the new model has a bound state, too. But with a binding energy of about 400 MeV below the $`\mathrm{\Lambda }N`$ threshold it is located completely outside of the physical region. One could speculate, of course, that this bound state is a manifestation of the Pauli forbidden $`(11)_s`$ state at the quark level FujiwaraP .) Furthermore, the cusp structure at the opening of the $`\mathrm{\Sigma }N`$ threshold is much less pronounced in the new model. In the old model this structure was primarily caused by a large amplitude in the tensor-coupled $`{}_{}{}^{3}S_{1}^{}^3D_1`$ partial wave of the $`\mathrm{\Lambda }N`$$`\mathrm{\Sigma }N`$ transition. This amplitude is now much smaller. As a consequence also the transition cross section for $`\mathrm{\Sigma }^{}p\mathrm{\Lambda }n`$ is now somewhat smaller, though still in line with the empirical informations. In the $`\mathrm{\Sigma }^{}p`$ channel the new model yields a stronger energy dependence of the reaction cross section as it is favoured by the available cross-section data. In the other two measured reaction channels the agreement with the data is equally good, if not better, for the new model.
Note that the $`\mathrm{\Sigma }^+p`$ and $`\mathrm{\Sigma }^{}p`$ elastic cross sections are not “true” total cross sections. The cross sections that were measured are defined as Eisele
$$\sigma =\frac{2}{\mathrm{cos}\theta _{\mathrm{max}}\mathrm{cos}\theta _{\mathrm{min}}}_{\mathrm{cos}\theta _{\mathrm{min}}}^{\mathrm{cos}\theta _{\mathrm{max}}}\frac{d\sigma (\theta )}{d\mathrm{cos}\theta }d\mathrm{cos}\theta ,$$
(16)
with typical values $`0.2`$ to $`0.5`$ for $`\mathrm{cos}\theta _{\mathrm{min}}`$ and $`0.3`$ to $`0.5`$ for $`\mathrm{cos}\theta _{\mathrm{max}}`$. In order to stay as close as possible to the plotted experimental data, the theoretical curves in Figs. 9(c) and (d) have been calculated with $`\mathrm{cos}\theta _{\mathrm{min}}=0.5`$ and $`\mathrm{cos}\theta _{\mathrm{max}}=0.5`$.
Cross sections at somewhat higher energies are presented in Fig. 8. Note that the data shown in this figure have not been taken into account in the fitting process and therefore the results are genuine predictions of the model. Also here the agreement with the data is satisfactory.
The differential $`YN`$ scattering cross sections presented in Fig. 9 are likewise genuine predictions of our $`YN`$ model. We want to point out that the empirical information in those figures comes from data taken from a finite momentum interval, e.g. 160 $`<p_{Lab}<`$ 180 MeV/c for the $`\mathrm{\Sigma }^+p`$ channel Alex , whereas the calculations were performed for the central value of that momentum interval as it is given in the various plots. Note also that the original $`YN`$ model of the Jülich group was fitted to the data without including the Coulomb interaction (and without taking into account the mass splitting between $`\mathrm{\Sigma }^{}`$, $`\mathrm{\Sigma }^0`$, and $`\mathrm{\Sigma }^+`$). Thus, the corresponding results presented in Fig. 9 do not show the strong forward peak caused by the Coulomb amplitude in the charged channels.
Evidently, also the data on differential cross sections are rather well reproduced by our new $`YN`$ model. In comparison to the results of the original Jülich model one can say that the angular dependence in the $`\mathrm{\Sigma }^{}p`$ channel is now much better described and it seems to be more in line with the trend of the angular dependence exhibited by the data in the $`\mathrm{\Sigma }^{}p\mathrm{\Lambda }n`$ channel too.
The dashed curves in Figs. 7, 8 and 9 are results from an alternative model where the contributions from the disputed $`\kappa `$(1000) meson have been replaced by a contact interaction. Obviously there is practically no sensitivity to the concrete range of the contribution in the scalar channel with isospin 1/2 – besides that it has to be of fairly short range. In this context we want to mention that we could achieve a comparable description of the data even with a $`\kappa `$ mass as low as 800 MeV Aitala .
For exploring the differences between the original Jülich $`YN`$ model and our new model in more detail we present in Figs. 10, 11 further observables where, however, no data are available. Fig. 10 contains differential cross sections, polarizations and the depolarization parameter $`D_{nn}`$ (definition and explicit expressions for those observables can be found in the appendix B of Ref. Reu ) for the $`\mathrm{\Lambda }N`$ channel. We present predictions at two energies, one ($`p_{lab}`$ = 150 MeV/c) close to the $`\mathrm{\Lambda }N`$ threshold and one ($`p_{lab}`$ = 600 MeV/c) close to (but below) the $`\mathrm{\Sigma }N`$ threshold. The results at the higher energy reveal that the new model differs drastically from the old one. The differential cross section in the new model is strongly forward-peaked whereas the one of the old models peaks in forward and backward direction. The polarization and $`D_{nn}`$ have even different signs. The observables at the lower energy are still dominated by the $`S`$ waves and therefore exhibit only minor differences. But one can see from the differential cross section that the onset of higher partial waves occurs earlier for the new $`YN`$ model.
Similarly striking differences are present also in the predictions for other differential observables though we refrain from showing them here.
For the various $`\mathrm{\Sigma }N`$ channels we present predictions for the polarization and the depolarization parameter $`D_{nn}`$ at $`p_{lab}`$ = 500 MeV/c, cf. Fig. 11. Also here one can see that, in general, there are large differences between the results of the old and the new model.
### IV.4 Low energy parameters and phase shifts
For the computation of the low energy parameters and phase shifts we omit the Coulomb interaction and ignore the mass differences between the $`\mathrm{\Sigma }`$’s and proton and neutron so that we can solve the Lippmann-Schwinger equation in isospin basis. This allows us to present also results for the $`\mathrm{\Sigma }N`$ system in the $`I=1/2`$ channel. The $`YN`$ S-wave low energy parameters are listed in Table 4 while phase shifts for selected partial waves are shown in Figs. 12 and 13.
From Table 4 one can see that the scattering lengths in the $`{}_{}{}^{3}S_{1}^{}`$ $`\mathrm{\Lambda }N`$ partial wave ($`a_t`$) are of similar magnitude for the old and new $`YN`$ models, but in the $`{}_{}{}^{1}S_{0}^{}`$ state ($`a_s`$) the new model yields a significantly larger value. The stronger $`{}_{}{}^{1}S_{0}^{}`$ component of the new model is reflected in the larger $`\mathrm{\Lambda }p`$ cross section near threshold, cf. Fig. 7, and it is expected to provide sufficient strength in order to support a bound hypertriton state Nogga . Indeed recent $`YN`$ models like NSC97f of the Nijmegen group NijV or the Ehime model 00A Ehime , that apparently lead to a bound hypertriton Miya3 ; Ehime , predict singlet scattering lengths that are very similar to that of our new model.
In this context we want to mention that the static version of the old Jülich $`YN`$ model Reu did not support a hypertriton bound state Miya1 . However, in that model both the $`{}_{}{}^{1}S_{0}^{}`$ as well as the $`{}_{}{}^{3}S_{1}^{}`$ $`\mathrm{\Lambda }N`$ scattering lengths are considerably smaller Reu than in our new $`YN`$ model.
The scattering lengths and effective ranges for $`\mathrm{\Sigma }N`$ with $`I=1/2`$ are complex because this channel is coupled to the $`\mathrm{\Lambda }N`$ system. In the singlet case the scattering lengths are comparable for the two models whereas in the triplet case they even have opposite signs. We want to emphasize, however, that in both models the latter partial wave is attractive. But in the original Jülich model the attraction is so strong that there is a near-threshold quasibound state in the $`\mathrm{\Sigma }N`$ channel that causes the real part of $`a_t`$ to be positive - like in case of the corresponding $`NN`$ partial wave and the deuteron. Let us mention that practically the same situation occurs in the Nijmegen model NSC97f, whose pole structure has been investigated and thoroughly discussed in Ref. Yamamura . As a consequence of the near-threshold pole both these models yield a very pronounced cusp-like structure in the $`\mathrm{\Lambda }p`$ cross section at the opening of the $`\mathrm{\Sigma }N`$ channel, cf. Fig. 7 and Fig. 2 in Ref. NijV , respectively. In our new $`YN`$ model, on the other hand, the cusp at the $`\mathrm{\Sigma }N`$ threshold is much less pronounced. Note that the $`{}_{}{}^{1}S_{0}^{}`$ partial wave is attractive too. As already mentioned above, in the original Jülich model there is a bound state in the $`\mathrm{\Sigma }N`$ channel – as evidenced by the broad bump in the $`\mathrm{\Lambda }p`$ cross section around $`p_{lab}`$ 350 MeV/c. And the new $`YN`$ model has also a bound state which is located, however, around 400 MeV below the $`\mathrm{\Lambda }N`$ threshold and therefore completely outside of the physically relevant region.
Let us finally come to the $`\mathrm{\Sigma }N`$ channel with $`I=3/2`$. Here we see that the singlet scattering length of the new model is about twice as large as the one of the original Jülich model. Note that a comparably large singlet scattering length is also predicted by all of the $`YN`$ models presented in Ref. NijV . The scattering lengths for $`{}_{}{}^{3}S_{1}^{}`$ are small in both cases, but of opposite sign. Now, however, it is indeed so that our new $`YN`$ model is repulsive in this partial wave whereas the old model is attractive. It is interesting that basically all available $`YN`$ models predict rather small values for the spin-triplet scattering length of the $`\mathrm{\Sigma }N`$ $`I=3/2`$ channel Holz ; NijIII ; NijIV ; NijV ; Fujiwara , though there is no general trend as far as the sign is concerned. We also observe an unnaturally large value for the triplet effective range, which is clearly related to the strong suppression of the corresponding scattering length. Such a scenario will require special attention when this channel is considered in effective field theory (for further discussion, see Sec. V).
Predictions for $`\mathrm{\Lambda }N`$ and $`\mathrm{\Sigma }N`$ phase shifts for selected $`S`$\- and $`D`$-waves are shown in Fig. 12 and those for $`P`$-waves can be found in Fig. 13. The $`\mathrm{\Sigma }N`$ $`S`$-wave phase shifts reflect the features that we already discussed in the context of the scattering lengths. For example one can see that the phase shift for the $`{}_{}{}^{3}S_{1}^{}`$ $`I=1/2`$ state starts at 180<sup>0</sup> for the original Jülich model, as it is expected for a partial wave where a bound state is present. For the $`I=3/2`$ state the corresponding phase is positive, reflecting an attractive interaction, whereas the phase shift resulting from the new $`YN`$ model is negative. Note that the phases for $`{}_{}{}^{1}S_{0}^{}`$ and $`I=1/2`$ should both start at 180<sup>0</sup> because, as mentioned above, there is a bound state in both models.
The opening of the $`\mathrm{\Sigma }N`$ channel at around $`E_{Lab}`$ 170 MeV is cleary reflected in the $`{}_{}{}^{3}S_{1}^{}`$ phase shift of the $`\mathrm{\Lambda }N`$ system. But its effect on the $`{}_{}{}^{3}D_{1}^{}`$ phase shift is even more striking where, for the old Jülich model, the phase even goes through 90 degrees. In fact, the resonance-like behaviour in that partial wave is predominantly responsible for the strong enhancement of the $`\mathrm{\Lambda }N`$ cross section in the vicinity of the $`\mathrm{\Sigma }N`$ threshold, cf. Fig. 7. In addition, the transition amplitude $`{}_{}{}^{3}D_{1}^{}(\mathrm{\Lambda }N)^3S_1(\mathrm{\Sigma }N)`$ provides a significant contribution to the $`\mathrm{\Sigma }^{}p\mathrm{\Lambda }n`$ cross section. In the new model the $`{}_{}{}^{3}D_{1}^{}`$ phase shift of the $`\mathrm{\Lambda }N`$ system is much smaller. Accordingly, the cusp-like structure at the $`\mathrm{\Sigma }N`$ threshold is much less pronounced and the $`\mathrm{\Sigma }^{}p\mathrm{\Lambda }n`$ cross section is somewhat reduced in this model, as can be seen in Fig. 7.
The predictions for the $`P`$ waves (Fig. 13) show a varying picture. In the $`\mathrm{\Lambda }N`$ system most of the phases are now attractive whereas they are mostly repulsive for the old model. This concerns in particular the $`{}_{}{}^{1}P_{1}^{}`$ amplitude, which is fairly large in the new model, but also the $`{}_{}{}^{3}P_{0}^{}`$ partial wave. In the $`I=3/2`$ channel of the $`\mathrm{\Sigma }N`$ system the results of the two models are qualitatively rather similar. To some extent this is also the case for the $`I=1/2`$ channel though here the $`{}_{}{}^{3}P_{1}^{}`$ amplitude of the new $`YN`$ model is significantly larger than the one of the old Jülich model. Indeed the simultaneous enhancement in the $`{}_{}{}^{3}P_{1}^{}`$ ($`\mathrm{\Sigma }N`$) and $`{}_{}{}^{1}P_{1}^{}`$ ($`\mathrm{\Lambda }N`$) phase shifts is caused by a stronger antisymmetric spin-orbit force between the $`\mathrm{\Lambda }N`$ and $`\mathrm{\Sigma }N`$ channels in the new model. The increase is primarily due to the $`\rho `$ exchange contribution whose strength for the $`\mathrm{\Lambda }N\mathrm{\Sigma }N`$ transition, fixed from correlated $`\pi \pi K\overline{K}`$ exchange, is about twice as large as what was used in the old Jülich model, cf. Table 11 of Ref. REUBER . In this context let us mention that some other $`YN`$ models exhibit a similarly strong coupling between those partial waves and channels FujiwaraA .
## V Summary and outlook
We have presented a meson-exchange model of the $`YN`$ interaction where – as the main new feature – the contributions both in the scalar-isoscalar ($`\sigma `$) and the vector-isovector ($`\rho `$) channels are constrained by a microscopic model of correlated $`\pi \pi `$ and $`K\overline{K}`$ exchange.
An essential part of baryon-baryon interactions is the strong medium-range attraction, which in one-boson-exchange models is parameterized by exchange of a fictitious scalar-isoscalar meson with mass around 500 MeV. In extended meson exchange models this part is naturally generated by two-pion exchange contributions. As well as uncorrelated two-pion exchange, correlated contributions must be included in which the exchanged pions interact during their exchange; these terms in fact provide the main contribution to the intermediate-range interaction.
As kaon exchange is an essential part of hyperon-nucleon interactions a simultaneous investigation of correlated $`\pi \pi `$ and $`K\overline{K}`$ exchanges is clearly necessary. In Ref. REUBER the correlated $`\pi \pi `$ and $`K\overline{K}`$ exchange contributions in various baryon-baryon channels have therefore been investigated within a microscopic model for the transition amplitudes of the baryon-antibaryon system ($`B\overline{B^{}}`$) into $`\pi \pi `$ and $`K\overline{K}`$ for energies below the $`B\overline{B^{}}`$ threshold. The correlations between the two mesons have been taken into account by means of $`\pi \pi K\overline{K}`$ amplitudes, determined in the field theoretical framework of Refs. Lohse ; Janssen ; Schutz , which provide an excellent description of empirical $`\pi \pi `$ data up to 1.3 GeV. With the help of unitarity and dispersion-theoretical methods, the baryon-baryon amplitudes for correlated $`\pi \pi `$ and $`K\overline{K}`$ exchange in the $`J^P=0^+`$ ($`\sigma `$) and $`J^P=1^{}`$ ($`\rho `$) $`t`$-channels have then been determined. With this model it is possible to reliably take into account correlated $`\pi \pi `$ and $`K\overline{K}`$ exchange in both the $`\sigma `$ and $`\rho `$ channels for various baryon-baryon reactions. Given the strong constraints on $`\sigma `$ as well as $`\rho `$ exchange from correlated $`\pi \pi `$ exchange, a more sound microscopic model for the $`YN`$ interaction can hence now be constructed.
Besides contributions from correlated $`\pi \pi `$ and $`K\overline{K}`$ exchange the present model incorporates also the standard one-boson exchanges of the lowest pseudoscalar and vector meson multiplets with coupling constants fixed by SU(6) symmetry relations. Thus, in the present model the long- and intermediate-range part of the $`YN`$ interaction is completely determined – either by SU(6) constraints or by correlated $`\pi \pi `$ and $`K\overline{K}`$ exchange.
In addition there are some short-ranged ingredients. First of all, the contribution from the $`a_0(980)`$ meson is taken into account. Secondly, we consider the exchange of a strange scalar meson, the $`\kappa `$, with mass $`1000`$ MeV. (Note that these pieces were not taken into account in the earlier $`YN`$ models of the Jülich group Holz ; Reu .) These short-ranged contributions are also viewed as being due to correlated meson-meson exchanges but in practice they are parametrized phenomelogically in terms of one-boson-exchange contributions in the corresponding spin-isospin channels. In particular, no SU(3) relations are imposed on the short-range part. This assumption is based on our observation that the contributions in the $`\rho `$ exchange channel as they result from correlated $`\pi \pi `$ and $`K\overline{K}`$ no longer fulfill SU(3) relations, but it also acknowledges the fact that at present there is no general agreement about who are the actual members of the lowest-lying scalar meson SU(3) multiplet.
The new $`YN`$ model provides a rather satisfactory reproduction of the available $`YN`$ data. It describes not only the integrated cross sections for $`\mathrm{\Lambda }p`$ and the various $`\mathrm{\Sigma }N`$ channels but also the few available data on differential cross sections, even though the latter were not included in the fitting procedure. We see that as an indication that the data are compatible with the assumption of SU(6) symmetry for the pseudoscalar sector of our $`YN`$ model.
As the main qualitative difference between the old $`YN`$ Jülich model Holz (with standard $`\sigma `$ and $`\rho `$ exchange) we want to mention that the broad shoulder at $`p_{lab}`$ 350 MeV/c in the $`\mathrm{\Lambda }p\mathrm{\Lambda }p`$ channel, predicted by that model but not seen in the experiments, is no longer present in the new model. But, as a more detailed comparison revealed, there are also striking differences between these two models in the predictions for the individual partial waves. For example, in the new model the triplet $`S`$ wave in the $`I=3/2`$ channel of the $`\mathrm{\Sigma }N`$ system is repulsive and some of the $`P`$-wave amplitudes are significantly larger. Thus, it will be interesting to see the performance of the new $`YN`$ interaction model in applications to few- and many-body systems involving hyperons Nogga .
This study also paves the way for a systematic investigation in the framework of effective field theory, see Henk . In such a framework, pion- and kaon exchange supplemented by four-baryon contact interactions (these encode the contributions from the exchange of heavier mesons not linked to chiral symmtery) is considered to generate a potential based on the power counting rules. It remains to be seen how well such a more systematic approach can indeed describe the data and what conclusions can be drawn about three-baryon forces that naturally arise in such a framework.
## Acknowledgements
We thank J. Speth and W. Melnitchouk for collaboration during the early stages of this investigation. We also thank A. Nogga for a careful reading of our manuscript. This research is part of the EU Integrated Infrastructure Initiative Hadron Physics Project under contract number RII3-CT-2004-506078. The work was supported in part by DFG through funds provided to the special research grant TR-16 “Subnuclear Structure of Matter”.
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# Comment on the ”Origin of the opalescence at the 𝜶↔𝜷 transition of quartz: role of the incommensurate phase studied by synchrotron radiation”
## Abstract
The interpretation of the origin of the light scattering anomalies near the transition from incommensurate (IC) to the $`\alpha `$-phase of quartz by Dolino et al. (Phys. Rev. Lett. 94, 155701 (2005)) is commented. It is shown that the observed IC structure is a pure transverse acoustic (TA) modulation without any soft optic mode component. Such a modulation cannot be responsible for the observed light scattering contrary to the interpretation by Dolino et al.
incommensurate; phase transition
The present paper is a comment on the interpretation by Dolino et al. 1 of their recent observations in quartz. In the paper1 , from the diffraction experiment using synchrotron radiation, the authors showed that the incommensurate (IC) modulation wavevector in quartz becomes extremely small in the region of the sample close to the boundary between the IC phase and the $`\alpha `$\- phase. In this region, where the two phases coexist in the same sample, very intense light scattering is observed.shustin Therefore, it is very important to consider how the intense light scattering is related to the reported extraordinary long period modulation.
In the following, we show that (i) the observed long-period IC modulation in quartz is a pure transversal acoustic (TA) modulation, including neither longitudinal acoustic (LA) component nor the optic mode component. (ii) A long-period acoustic modulation cannot appear spontaneously as a result of the IC phase transition associated with the acoustic mode softening, and the mechanism of its formation should be understood. (iii) Such a modulation cannot be at the origin of any observed light scattering anomalies contrary to the interpretation by Dolino et al.1 .
The consideration below is mainly based on the observations by Dolino et al. 3 of the IC satellite reflections in the elastic neutron diffraction, and their expected evolving in view of the recent observations in the synchrotron radiation 1 . As observed (see Fig. 1), near the Bragg reflections (100), (200) …the satellites, corresponding to the IC vector $`𝐤`$ parallel to the main Bragg reflection vector (l00), are systematically absent, i.e. four satellites are observed instead of six. Similar observations by Gouhara and Kato 5 ; 6 brought them to a conclusion, that the IC modulation in quartz is mainly of an acoustic character, with a small (of about 0.1 fraction) optical atomic displacements component. The analysis carried out by the present authors 7 , based on the comparison of the IC satellites near the (l00) and (l0m)-type Bragg reflections, showed that the IC modulation in quartz is a pure TA modulation, including, as one of its components, the transversal $`u_z`$-displacements. The contribution of the TA $`u_z`$-displacements to the diffraction was interpreted by Gouhara and Kato 5 ; 6 as that from a small optical component. The existence of $`u_z`$-displacements was also demonstrated by the recent MD calculations. Dmitirev Below we bring new arguments, unambiguously demonstrating the TA character of the IC modulation in quartz. However, contribution of a long-period acoustic modulation to the dielectric constant’s components (i.e. to the light scattering) is proportional to the square of the IC wavevector, and therefore it must be negligibly small for the small IC modulation vectors observed in quartz ($`k<0.03a^{}`$ and decreases on cooling, according to the recent observations1 , down to $`0.002a^{}`$).
We show that in the case if the IC modulation contains an optic mode displacements component, an intense IC satellite reflections should necessarily appear in the positions of the missing satellites (in paricular, near (300) in Fig. 1).
Let us assume that the IC modulation in quartz, however, contains an optical component. In such a case, the long-period triple-$`k`$ IC modulation of the quartz’s $`\alpha \beta `$ transition parameter $`\eta `$ necessarily should induce a longitudinal acoustic (LA) modulation with a large amplitude $`u_{l0}`$. The amplitude $`u_{l0}`$ of such LA modulation should be of the same order of magnitude as that for the $`\eta `$\- modulation for the IC vectors $`k0.03a^{}`$, and for the IC vectors $`k0.002a^{}`$ (near the transition to the $`\alpha `$-phase) it should increase up to $`10`$ times (the corresponding diffraction satellites should grow in intensity up to $`10^2`$ times). As a result, six intense IC satellites should be induced by the LA modulation instead of the four satellites near any (l00) reflection.
The elastic strain dependent part of the quartz’s thermodynamic potential is of the form 4 :
$`{\displaystyle }\{a[(u_{xx}u_{yy}){\displaystyle \frac{\eta }{x}}2u_{xy}{\displaystyle \frac{\eta }{y}}]+r\eta ^2(u_{xx}+u_{yy})+`$
$`{\displaystyle \frac{c_{11}c_{66}}{2}}(u_{xx}+u_{yy})^2+{\displaystyle \frac{c_{66}}{2}}[(u_{xx}u_{yy})^2+4u_{xy}^2]\}dV`$ (1)
where $`u_{xy},u_{xx},u_{yy}`$ are the elastic strains, $`a`$ and $`r`$ are some coefficients of expansion, which are of an atomic order of magnitude. For the triple-$`k`$ IC optical modulation
$$\eta (R)=\eta _0\mathrm{cos}(\mathrm{𝐤𝐑}+\phi )+\eta _0\mathrm{cos}(𝐤_\mathrm{𝟏}𝐑+\phi _1)+\eta _0\mathrm{cos}(𝐤_\mathrm{𝟐}𝐑+\phi _2)$$
with $`𝐤+𝐤_1+𝐤_2=0`$, and for the acoustic displacements vector $`𝐮=𝐮_l+𝐮_t`$, presented as the sum of the longitudinal $`𝐮_l`$ (along the IC vector $`k`$) and transversal $`𝐮_t`$ displacements, one can minimize Eq. (1) with respect to $`𝐮_l`$, and obtain an equilibrium LA modulation wave:
$$𝐮_l=𝐮_{l0}\mathrm{sin}(\mathrm{𝐤𝐑}+\phi _1+\phi _2)+𝐮_{l0}^{}\mathrm{cos}(\mathrm{𝐤𝐑}+\phi ),$$
with the amplitudes
$$u_{l0}=\frac{2r\eta _0^2}{c_{11}k}\mathrm{and}u_{l0}^{}=\frac{a\eta _0\mathrm{cos}3\varphi }{c_{11}},$$
where $`\varphi `$ is the angle between the $`x`$ axis and the IC vector $`𝐤`$. In the observed IC structure the angle $`\varphi `$ is close to $`\pi /6`$, and therefore the corresponding amplitude $`u_{l0}^{}`$ is close to zero and generally neglected. An identical LA IC amplitudes must exist for all the six wavevectors $`\pm 𝐤,\pm 𝐤_1`$ and $`\pm 𝐤_2`$.
Taking into account that the IC wavevector in quartz is very small, and it decreases on cooling from 0.03$`a^{}`$ (as in Fig. 1) down to 0.002$`a^{}`$ (according to the recent publication by Dolino1 ), one can see the enormous increase of LA modulation’s amplitude $`u_{l0}`$ (see the $`k`$ vector in the denominator).
The physical reason of the large LA amplitude is very plain. Formation of a long-period acoustic modulation requires a very small energy (this energy is proportional to $`k^2`$, and is tending to zero with $`k0`$). Though the acoustic-optic coupling term $`\eta ^2u_{ii}`$ in Eq. (1) is rather small, however it is sufficient for inducing an enormous LA amplitude for the small $`k`$-s.
For calculation of the X-ray diffraction scattering amplitude, induced by the IC wave, one should expand the crystals density function as:
$`F_G\mathrm{exp}[i𝐆(𝐑+𝐮)]F_G\mathrm{exp}[i\mathrm{𝐆𝐑}]+`$
$`iF_G(\mathrm{𝐆𝐮}_0)\mathrm{exp}[i(𝐆\pm 𝐤)𝐑],`$ (2)
where $`F_G`$ is the structure factor, corresponding to the $`𝐆`$-Bragg reflection, $`𝐮_0`$ is the amplitude of the IC acoustic modulation $`𝐮(𝐑)`$, and the scalar product $`(\mathrm{𝐆𝐮}_0)=u_{l0}G`$ for the case, when vector $`𝐆`$ is parallel to the IC vector $`𝐤`$ (as for the case of the missing satellites in Fig. 1). The first term in Eq. (2) gives diffraction to the main Bragg reflection $`𝐆`$ with intensity $`|F_G|^2`$. The second term gives the acoustic modulation’s contribution to the satellite reflections $`𝐆\pm 𝐤`$ with intensity $`|F_G(\mathrm{𝐆𝐮}_0)|^2`$. For the case of $`𝐆𝐤`$, this intensity takes the form $`|F_GGu_{l0}|^2`$, and it extremely increases with $`k0`$, since, as shown above, $`u_{l0}`$ increases as $`1/k`$.
For the case $`k=0.03a^{}`$, it is easy to estimate that the satellite intensity induced by $`u_{l0}`$ and contributing to the position of the missing satellites in Fig. 1 should be of the same order of magnitude as any intensity, induced by the optical $`\eta `$-atomic displacements. Note that the satellite intensity contributed from the optical component is $`|F_GG\eta _0|^2`$. In the vicinity of $`1K`$ of the phase transition, the optical amplitude $`\eta _010^2\eta _{at}`$, where $`\eta _{at}`$ is of an atomic order of magnitude, $`k10^2a^{}`$, and subsequently, $`u_{l0}`$ and $`\eta _0`$ are of the same order of magnitude. For the smaller IC vectors $`k0.002a^{}`$, $`u_{l0}`$ is larger than $`\eta _0`$ about $`10`$ times, and the ratio of the corresponding intensities should make $`10^2`$. In other words, for $`k0.03a^{}`$ six reflections of about the same intensity in Fig. 1 should be observed. And besides, since the scalar product $`(\mathrm{𝐆𝐮}_0)`$ in Eq. (2) is maximal for the case when $`𝐆`$ and $`𝐤`$ are parallel, the intensity of the $`𝐆\pm 𝐤`$ satellites should be larger, than the intensity of the $`𝐆\pm 𝐤_1`$ and $`𝐆\pm 𝐤_2`$ satellites. In other words, the missing satellites in Fig. 1 should be detected as the most intensive, compared with four other satellites near (l00).
The only explanation of the systematic observation of four satellites near the (l00) type Bragg reflections is the TA character of the IC modulation, without optical and LA components, since otherwise, any optical component in the triple-k modulation necessarily induces a strong LA component, which will contribute to the positions of the missing satellites.
The TA IC structure, which, in fact, exists in quartz, cannot noticeably contribute to the light scattering anomalies, since its effect on the dielectric constant $`\epsilon _{ij}`$ is proportional to $`u_{xy}^2`$ (i.e. to a very small parameter $`k^2`$), while the light scattering anomalies are observed just when $`k0`$. So, we do not agree with the interpretation of the light scattering anomalies in quartz by Dolino et al. 1 , attributed to the observed IC modulation. Neither the domains of the IC structure (the adjacent domains with different values of $`k`$ in the small-angle scattering zone 1 ), nor the Dauphine twins observed in the ”fog zone” can be responsible for the light scattering in quartz. The latter was discussed in detail and estimated by the present authors earlier 8 , where the origin of the intense light scattering was attributed to the ferroelastic domains. The arguments given above completely contradict the observations and reject the traditional model for the IC transition adopted by Dolino et al. 1 .
Finally we introduce another argument, supporting the above discussion. As it follows from the symmetry, in case of the triple-$`k`$ IC optical modulation, the dielectric constant’s components $`\mathrm{\Delta }\epsilon _{xx},\mathrm{\Delta }\epsilon _{yy},\mathrm{\Delta }\epsilon _{zz}`$ are proportional to $`\eta _0^2\mathrm{sin}(\mathrm{𝐤𝐑}+\phi _1+\phi _2)`$. Since the magnitude of the IC vector $`𝐤`$ decreases down to 0.002$`a^{}`$, then the light of 500nm or smaller wavelength (i.e. with wavevector $`0.001a^{}`$), should diffract in back-direction (or close to it) in such IC modulation. Such a feature, which can be visually detected, is not observed in quartz. This fact also proves, that the LA modulation (and subsequently an optical modulation), which directly follows from the traditional model used by Dolino et al. 1 , does not exist in the IC phase of quartz.
It should be mentioned that a long-period acoustic modulation cannot appear spontaneously, as a result of the IC phase transition, associated with the acoustic mode softening. Such a mode softening would have a lot of consequences, including a significant drop in the elastic constant $`c_{66}`$, which have never been observed in quartz.kimizuka One may assume that some latent (driving) phase transition takes place in quartz near the $`\alpha \beta `$ transition, and the observed acoustic modulation is only a manifestation of this hidden transition. However, a pure TA character of the IC modulation in quartz was sufficiently simply and reliably demonstrated also in the earlier consideration 7 , based on the different arguments. So, it cannot be ignored, and understanding of its origin is of a priority importance in the quartz’s phase transition problem.
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# Linear relations among holomorphic quadratic differentials and induced Siegel’s metric on ℳ_𝑔
## 1. Introduction
In spite of the remarkable progresses in understanding the Schottky problem, the characterization of the Schottky locus as the zero set of modular forms on the Siegel’s upper half-space remains a fundamental open question. Such a question is strictly related to the problem of characterizing the Schottky locus by means of $`(g2)(g3)/2`$ linearly independent conditions. As suggested also by Mumford (see pg. 241 of ), a possible solution of that problem should follow by a deeper understanding of Petri’s construction . Actually, since Petri’s theorem determines the ideal of canonical curves of genus $`g4`$ by means of linear relations among holomorphic abelian differentials, it seems the natural framework for such an investigation.
Let $`\{\eta _1,\mathrm{},\eta _g\}`$ be the Petri’s basis for $`H^0(K_C)`$, with $`C`$ a canonical curve of genus $`g`$. In Petri’s work the coefficients $`C_{ij}^k`$ in the relationships among quadratic differentials $`_{i,j}^gC_{ij}^k\eta _i\eta _j=0`$, $`k=1,\mathrm{},(g2)(g3)/2`$, are not determined. Finding such coefficients is a necessary condition for an explicit characterization of the ideal of canonical curves. Here, we express Petri’s relation in determinantal form, so that, besides the explicit determination of the coefficients, it is shown that the locus of canonical curves corresponds to a determinantal variety.
We introduce modular invariant bases for holomorphic differentials, leading to a refinement of Petri’s basis and to an immediate derivation of Fay’s trisecant identity . A key point is the introduction of a indexing, which includes the combinatorics of the Petri construction, mapping the components of matrices in the Siegel upper half-space to vector components. This provides the volume form on the moduli space $`\widehat{}_g`$ of canonical curves induced by the Siegel metric which, remarkably, is expressed in terms of the period Riemann matrix only. By the Kodaira-Spencer map, the above relations lead to an expression of the metric on $`\widehat{}_g`$, induced by the Siegel metric, that corresponds to the square of the Bergman reproducing kernel.
In the case of branched covering of the torus, corresponding to Jacobians with a distinguished complex multiplication , the derived relations should lead to identities of number theoretical interest. Our results, of interest also in superstring theory , provide the key for the $`(g2)(g3)/2`$ combinatorial $`\theta `$-identities in .
## 2. Determinantal characterization of canonical curves
Let $`C`$ be a canonical curve of genus $`g4`$ and $`\{\omega _i\}_{iI_g}`$, $`I_n:=\{1,\mathrm{},n\}`$, a basis of $`H^0(K_C)`$, with $`K_C`$ the canonical line bundle of $`C`$. Denote by $`\widehat{}_g`$ the corresponding locus in the moduli space $`_g`$ of compact Riemann surfaces. Each element of $`H^0(K_C^2)`$ can be written as a linear combination of the $`M:=g(g+1)/2`$ elements in
$$𝒮:=\{\omega _i\omega _j|ijI_g\}.$$
Since $`N:=h^0(K_C^2)=3g3`$, there are $`MN=(g2)(g3)/2`$ linearly independent relations among the quadratic differentials $`\omega _i\omega _j`$.
Let $`p_1,\mathrm{},p_g`$ and $`q_1,\mathrm{},q_{2g2}`$ be two sets of arbitrary points on $`C`$. Choose a local trivialization of the canonical line bundle and set
$$a_{ij,r}:=det\omega (p_1,\mathrm{},p_{i1},q_r,p_{i+1},\mathrm{},p_g)det\omega (p_1,\mathrm{},p_{j1},q_r,p_{j+1},\mathrm{},p_g),$$
where $`det\omega (x_1,\mathrm{},x_g):=det\omega _i(x_j)`$. Set
$$A(k,l):=\left(\begin{array}{ccccccc}a_{12,1}& \mathrm{}& a_{1g,1}& a_{23,1}& \mathrm{}& a_{2g,1}& a_{kl,1}\\ a_{12,2}& \mathrm{}& a_{1g,2}& a_{23,2}& \mathrm{}& a_{2g,2}& a_{kl,2}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ a_{12,2g2}& \mathrm{}& a_{1g,2g2}& a_{23,2g2}& \mathrm{}& a_{2g,2g2}& a_{kl,2g2}\end{array}\right),$$
$`3k<lg`$, $`g4`$.
###### Theorem 2.1.
$$detA(k,l)=0,3k<lg.$$
Set
$$\mathrm{\Delta }_{mn}:=()^{m+n}\underset{_{jn}^{im}}{det}A(k,l)_{ij},D_{pq}:=()^{p+q}\underset{_{jq}^{ip}}{det}\omega _j(p_i),$$
and denote by $`A_{ij,r}(k,l)`$, $`i,jI_g`$, $`rI_{2g2}`$, the matrix obtained from $`A(k,l)`$ by replacing the row $`(a_{12,r},\mathrm{},a_{kl,r})`$ with $`(D_{1i}D_{2j},\mathrm{},D_{ki}D_{lj})`$.
###### Corollary 2.2.
For each $`rI_{2g2}`$, the following relations
$$\underset{i,j=1}{\overset{g}{}}\frac{detA_{ij,r}(k,l)}{\mathrm{\Delta }_{r\mathrm{\hspace{0.17em}2}g2}}\omega _i\omega _j=0,$$
$`3k<lg`$, provide $`(g2)(g3)/2`$ linearly independent conditions on $`𝒮`$ which are independent of the points $`q_i`$, $`iI_{2g2}`$.
## 3. Distinguished bases of $`H^0(K_C^n)`$
Set $`N_n:=(2n1)(g1)+\delta _{1n}`$, $`n1`$, with $`\delta _{ij}`$ the Kronecker delta. Note that $`N_1=g`$ and $`N_2N`$. Fix a system of local coordinates on $`C`$.
###### Proposition 3.1.
Fix $`n1`$ and let $`p_1,\mathrm{},p_{N_n}`$ be a set of points in $`C`$ such that
$$det\varphi (p_1,\mathrm{},p_{N_n})0,$$
for an arbitrary basis $`\{\varphi _i\}_{iI_{N_n}}`$ of $`H^0(K_C^n)`$. Then
(3.1)
$$\gamma _i^n(z):=\frac{det\varphi (p_1,\mathrm{},p_{i1},z,p_{i+1},\mathrm{},p_{N_n})}{det\varphi (p_1,\mathrm{},p_{N_n})},$$
$`iI_{N_n}`$, determines a basis of $`H^0(K_C^n)`$ which is independent of the choice of the basis $`\{\varphi _i\}_{iI_{N_n}}`$ and, up to normalization, of the local coordinates on $`C`$.
###### Proof.
The matrix $`[\varphi ]_{ij}:=\varphi _i(p_j)`$ is non-singular. Then $`\gamma _i^n(z)=_j[\varphi ]_{ij}^1\varphi _j(z)`$, $`iI_{N_n}`$, is a basis of $`H^0(K_C^n)`$.∎
Note that $`\gamma _i^n(p_j)=\delta _{ij}`$, furthermore
(3.2)
$$det\gamma ^n(p_1,\mathrm{},p_{j1},z,p_{j+1},\mathrm{},p_{N_n})=\gamma _j^n(z).$$
###### Remark 1.
As we will see, the Fay trisecant identity directly follows by expressing Eq.(3.1) in terms of theta functions.
For $`n=1`$, the choice of $`g`$ points $`p_1,\mathrm{},p_gC`$, with $`det\omega _i(p_j)0`$, determines the basis $`\{\sigma _i\}_{iI_g}`$ of $`H^0(K_C)`$, where
(3.3)
$$\sigma _i(z):=\gamma _i^1(z),iI_g.$$
We now introduce a refinement of Petri’s basis for $`H^0(K_C^2)`$ which provides a modular invariant construction. Let us assume that the points $`p_1,\mathrm{},p_g`$ are in “general position” and that the effective divisor $`(\sigma _1)+(\sigma _2)_{i=3}^gp_i`$ consists of $`3g3`$ distinct points. Consider the following $`M`$ elements of $`H^0(K_C^2)`$
$$v_i:=\{\begin{array}{cc}\sigma _i^2,\hfill & iI_g,\hfill \\ & \\ \sigma _{j+k}\sigma _j,\hfill & i=k+j(2gj+1)/2,jI_{g1},kI_{gj}.\hfill \end{array}$$
###### Proposition 3.2.
$`\{v_i\}_{iI_N}`$ is a basis of $`H^0(K_C^2)`$.
###### Proof.
Set $`D:=_{i=3}^gp_i`$ and let us consider the effective divisors $`D_i:=(\sigma _i)D`$, $`i=1,2`$. Let us first prove that $`\sigma _i`$ is the unique $`1`$-differential, up to normalization, vanishing at $`D_i`$, $`i=1,2`$. Any $`1`$-differential $`\sigma _i^{}H^0(K_C)`$ vanishing at $`D_i`$, corresponds to an element $`\sigma _i^{}/\sigma _i`$ of $`H^0(𝒪(D))`$, the space of meromorphic functions $`f`$ on $`C`$ such that $`(f)+D`$ is an effective divisor. Suppose that there exists $`\sigma _i^{}`$ such that $`\sigma _i^{}/\sigma _i`$ is not a constant, so that $`h^0(𝒪(D))2`$. By Riemann-Roch theorem
$$h^0(K_C𝒪(D))=h^0(𝒪(D))\mathrm{deg}D1+g3,$$
so that there exist at least $`3`$ linearly independent $`1`$-differentials vanishing at $`D`$ and, in particular, there exists a linear combination of such differentials vanishing at $`p_1,\mathrm{},p_g`$. This implies that $`det\eta (p_1,\mathrm{},p_g)=0`$ for an arbitrary basis $`\{\eta _i\}_{iI_g}`$ of $`H^0(K_C)`$, contradicting the hypotheses. Fix $`\zeta _i,\zeta _{1i},\zeta _{2i}`$ such that
$$\underset{i=3}{\overset{g}{}}\zeta _i\sigma _i^2+\underset{i=1}{\overset{g}{}}\zeta _{2i}\sigma _1\sigma _i+\underset{i=2}{\overset{g}{}}\zeta _{1i}\sigma _2\sigma _i=0.$$
Evaluating this relation at the point $`p_j`$, $`3jg`$, yields $`\zeta _j=0`$. Set
(3.4)
$$t_1:=\underset{j=2}{\overset{g}{}}\zeta _{1j}\sigma _j,t_2:=\underset{j=1}{\overset{g}{}}\zeta _{2j}\sigma _j,$$
so that $`\sigma _1t_2=\sigma _2t_1`$. Since $`D`$, $`D_1`$ and $`D_2`$ consist of pairwise distinct points, $`t_i`$ vanishes at $`D_i`$, $`i=1,2`$ and then $`t_1/\sigma _1=t_2/\sigma _2=\zeta `$. By (3.4)
$$\zeta \sigma _1+\underset{j=2}{\overset{g}{}}\zeta _{1j}\sigma _j=0,\zeta \sigma _2\underset{k=1}{\overset{g}{}}\zeta _{2k}\sigma _k=0,$$
and, by linear independence of $`\sigma _1,\mathrm{},\sigma _g`$, we have $`\zeta =\zeta _{1j}=\zeta _{2k}=0`$, $`2jg`$, $`kI_g`$.∎
## 4. Proofs of Theorem 2.1 and Corollary 2.2
Let $`W(P)`$ be the Wronskian $`W(v_1,\mathrm{},v_N)(P)`$ of the basis $`\{v_i\}_{iI_N}`$ at a generic point $`PC`$, and $`\widehat{W}_{ij}(P):=W(v_1,\mathrm{},v_{i1},v_j,v_{i+1},\mathrm{},v_N)(P)`$.
###### Lemma 4.1.
The $`(g2)(g3)/2`$ linearly independent relations
(4.1)
$$v_i(z)W(P)=\underset{j=1}{\overset{N}{}}v_j(z)\widehat{W}_{ji}(P),$$
$`i=N+1,\mathrm{},M`$, hold $`zC`$.
###### Proof.
Immediate consequence of the Cramer rule. ∎
###### Remark 2.
The ratio $`\widehat{W}_{ij}(P)/W(P)`$ does not depend on $`P`$.
###### Remark 3.
Since for $`iI_g`$
(4.2)
$$\{\begin{array}{cc}v_j(p_i)=\delta _{ji},\hfill & jI_g,\hfill \\ & \\ v_j(p_i)=0,\hfill & j=g+1,\mathrm{},M,\hfill \end{array}$$
it follows that for $`z=p_i`$ Eq.(4.1) gives $`\widehat{W}_{ij}(P)=0`$ for $`iI_g`$ and $`j=N+1,\mathrm{},M`$.
###### Remark 4.
The Wronskians in the expansion (4.1) can be replaced by the corresponding determinant $`detv_j(x_i)`$, where $`x_1,\mathrm{},x_{3g3}`$ are arbitrary points on $`C`$.
Proof of Theorem 2.1. Assume that $`det\omega (p_1,\mathrm{},p_g)0`$. Define
$$x_i:=\{\begin{array}{cc}p_i,\hfill & iI_g,\hfill \\ & \\ q_{ig},\hfill & i=g+1,\mathrm{},N+1.\hfill \end{array}$$
Fix $`k,l`$ with $`3k<lg`$ and consider the matrix
$$\left(\begin{array}{cccc}v_1(x_1)& \mathrm{}& v_N(x_1)& \sigma _k(x_1)\sigma _l(x_1)\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ v_1(x_{N+1})& \mathrm{}& v_N(x_{N+1})& \sigma _k(x_{N+1})\sigma _l(x_{N+1})\end{array}\right).$$
By (4.2), this matrix has $`\mathrm{diag}(1,\mathrm{},1)`$ in the $`g\times g`$ upper left corner, $`0`$ in the $`g\times (2g2)`$ upper right corner and $`(det\omega (p_1,\mathrm{},p_g))^2A(k,l)`$ in the $`(2g2)\times (2g2)`$ lower right corner. On the other hand, by Lemma 4.1 and by Remark 4 the determinant of this matrix vanishes. Since such relations hold for $`(p_1,\mathrm{},p_g)`$ in a dense subset of $`C^g`$, they hold $`(p_1,\mathrm{},p_g)C^g`$ and the theorem follows.∎
Proof of Corollary 2.2. Divide the relations in Theorem 2.1 by $`\mathrm{\Delta }_{i\mathrm{\hspace{0.17em}2}g2}`$ and note that
$$a_{mn,r}=\underset{i,j=1}{\overset{g}{}}D_{mi}D_{nj}\omega _i(q_r)\omega _j(q_r).$$
Independence of the points $`q_i`$, $`iI_{2g2}`$, follows by noting that the coefficients of $`\omega _i\omega _j`$ in the relations are functions of $`q_i`$ with no zeroes or poles.∎
Define
$$(\mathfrak{1}_i,\mathfrak{2}_i):=\{\begin{array}{cc}(i,i),\hfill & 1ig,\hfill \\ (1,ig+1),\hfill & g+1i2g1,\hfill \\ (2,i2g+3),\hfill & 2gi3g3,\hfill \\ \mathrm{}\hfill & \mathrm{}\hfill \\ (g1,g),\hfill & i=g(g+1)/2,\hfill \end{array}$$
so that $`\mathfrak{1}_i\mathfrak{2}_i`$ is the $`i`$-th element in the $`M`$-tuple $`(11,22,\mathrm{},gg,12,\mathrm{},1g,23,\mathrm{})`$. $`u^g`$ and for all the $`g\times g`$ matrices $`A`$, set
(4.3)
$$uu_i:=u_{\mathfrak{1}_i}u_{\mathfrak{2}_i},(AA)_{ij}:=\frac{A_{\mathfrak{1}_i\mathfrak{1}_j}A_{\mathfrak{2}_i\mathfrak{2}_j}+A_{\mathfrak{1}_i\mathfrak{2}_j}A_{\mathfrak{2}_i\mathfrak{1}_j}}{1+\delta _{\mathfrak{1}_j\mathfrak{2}_j}},A_i:=A_{\mathfrak{1}_i\mathfrak{2}_i},$$
$`i,jI_M`$. In the following we will repeatedly make use of the identities
(4.4)
$$\underset{i,j=1}{\overset{g}{}}f(i,j)=\underset{ij}{\overset{g}{}}\frac{f(i,j)+f(j,i)}{1+\delta _{ij}}=\underset{k=1}{\overset{M}{}}\frac{f(\mathfrak{1}_k,\mathfrak{2}_k)+f(\mathfrak{2}_k,\mathfrak{1}_k)}{1+\delta _{\mathfrak{1}_k\mathfrak{2}_k}}.$$
In particular, if $`f(i,j)=f(j,i)`$, then
(4.5)
$$\underset{i,j=1}{\overset{g}{}}f(i,j)=\underset{k=1}{\overset{M}{}}(2\delta _{\mathfrak{1}_k\mathfrak{2}_k})f(\mathfrak{1}_k,\mathfrak{2}_k),$$
where we used the identity
$$2\delta _{ij}=\frac{2}{1+\delta _{ij}}.$$
With this notation, and observing that $`\sigma _i=_{j=1}^g[\omega ]_{ij}^1\omega _j`$, we have
(4.6)
$$v_i=\sigma \sigma _i=\underset{j=1}{\overset{M}{}}([\omega ]^1[\omega ]^1)_{ij}\omega \omega _j,iI_M.$$
Set $`w_{ij}:=W_{ij}/W`$, where $`W_{ij}(P):=W(v_1,\mathrm{},v_{i1},\omega \omega _j,v_{i+1},\mathrm{},v_N)(P)`$, and note that
(4.7)
$$\omega \omega _i=\underset{j=1}{\overset{N}{}}w_{ji}v_j,iI_M.$$
## 5. Siegel’s induced measure on $`\widehat{}_g`$ and Bergman reproducing kernel
Let
$$_g:=\{ZM_g()|{}_{}{}^{t}Z=Z,Y>0\},Y:=\mathrm{}Z,$$
be the Siegel upper half-space and $`\{\alpha _1,\mathrm{},\alpha _g,\beta _1,\mathrm{},\beta _g\}`$ a symplectic basis of $`H_1(C,)`$. Denote by $`\{\omega _i\}_{iI_g}`$ the basis of $`H^0(K_C)`$, dual of $`H_1(C,)`$, so that $`_{\alpha _i}\omega _j=\delta _{ij}`$, $`i,jI_g`$. Let $`\tau _{ij}:=_{\beta _i}\omega _j_g`$ be the Riemann period matrix of $`C`$. Under the symplectic transformation
$$\left(\begin{array}{c}\stackrel{~}{\alpha }\\ \stackrel{~}{\beta }\end{array}\right)=\left(\begin{array}{cc}D& C\\ B& A\end{array}\right)\left(\begin{array}{c}\alpha \\ \beta \end{array}\right),\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)Sp(2g,),$$
we have $`\stackrel{~}{\omega }={}_{}{}^{t}(C\tau +D)\omega `$, with $`\stackrel{~}{\tau }_{ij}`$ and $`\tau _{ij}`$ related by the modular transformation
$$\stackrel{~}{\tau }=(A\tau +B)(C\tau +D)^1.$$
Note that the basis $`\{\sigma _i\}_{iI_g}`$, defined in Eq.(3.3), is independent of the choice of the basis of $`H_0(K_C)`$ and therefore is modular invariant.
The Siegel metric
(5.1)
$$ds^2:=\mathrm{Tr}(Y^1dZY^1d\overline{Z}),$$
defines the volume form
$$d\nu =\frac{i^M}{2^g}\frac{_{ij}^g(\mathrm{d}Z_{ij}\mathrm{d}\overline{Z}_{ij})}{detY^{g+1}}.$$
We use the indexing introduced in Eq.(4.3) to express the Siegel metric on $`_g`$ where now the matrix elements $`Z_{ij}`$, $`i,jI_M`$, are seen as the components of the $`M`$-dimensional vectors $`Z:=(Z_1,\mathrm{},Z_M)`$.
###### Proposition 5.1.
(5.2)
$$ds^2=\underset{i,j=1}{\overset{M}{}}g_{ij}dZ_id\overline{Z}_j,$$
where
(5.3)
$$g_{ij}:=(2\delta _{\mathfrak{1}_i\mathfrak{2}_i})(Y^1Y^1)_{ij}.$$
###### Proof.
By (4.4) and (4.5)
$`ds^2`$ $`={\displaystyle \underset{i,j,k,l=1}{\overset{g}{}}}Y_{ij}^1dZ_{jk}Y_{kl}^1d\overline{Z}_{li}`$
$`={\displaystyle \underset{i,l=1}{\overset{g}{}}}d\overline{Z}_{li}{\displaystyle \underset{m=1}{\overset{M}{}}}{\displaystyle \frac{Y_{i\mathfrak{1}_m}^1Y_{l\mathfrak{2}_m}^1+Y_{i\mathfrak{2}_m}^1Y_{l\mathfrak{1}_m}^1}{1+\delta _{\mathfrak{1}_m\mathfrak{2}_m}}}dZ_{\mathfrak{1}_m\mathfrak{2}_m}`$
$`={\displaystyle \underset{m,n=1}{\overset{M}{}}}(2\delta _{\mathfrak{1}_n\mathfrak{2}_n})d\overline{Z}_{\mathfrak{1}_n\mathfrak{2}_n}{\displaystyle \frac{Y_{\mathfrak{1}_n\mathfrak{1}_m}^1Y_{\mathfrak{2}_n\mathfrak{2}_m}^1+Y_{\mathfrak{1}_n\mathfrak{2}_m}^1Y_{\mathfrak{2}_n\mathfrak{1}_m}^1}{1+\delta _{\mathfrak{1}_m\mathfrak{2}_m}}}dZ_{\mathfrak{1}_m\mathfrak{2}_m}`$
$`={\displaystyle \underset{m,n=1}{\overset{M}{}}}(2\delta _{\mathfrak{1}_n\mathfrak{2}_n})(Y^1Y^1)_{nm}dZ_md\overline{Z}_n.`$
Let $`k`$ be the Kodaira-Spencer map identifying the quadratic differentials on $`C`$ with the fiber of the cotangent of the Teichmüller space at $`C`$. We have
$$k(\omega _i\omega _j)=(2\pi i)^1d\tau _{ij}.$$
By Corollary 2.2 it follows that
(5.4)
$$\underset{i,j=1}{\overset{g}{}}\frac{detA_{ij,r}(k,l)}{\mathrm{\Delta }_{r\mathrm{\hspace{0.17em}2}g2}}d\tau _{ij}=0,$$
$`3k<lg`$. Set $`d\tau _i:=d\tau _{\mathfrak{1}_i\mathfrak{2}_i}`$, $`iI_M`$. Eq.(5.3) yields an explicit expression for the volume form on $`\widehat{}_g_g/Sp(2g,)`$ induced by the modular invariant Siegel metric on $`_g`$. Set $`\tau _2:=\mathrm{}\tau `$, and let $`\left|\tau _2^1\tau _2^1\right|_{j_1\mathrm{}j_N}^{i_1\mathrm{}i_N}`$, with $`i_k,j_k`$, $`kI_N`$, distinct elements of $`I_M`$, be the determinant of the $`N\times N`$ submatrix of $`(\tau _2^1\tau _2^1)_{ij}`$, built by taking the rows $`i_1,\mathrm{},i_N`$ and the columns $`j_1,\mathrm{},j_N`$.
###### Theorem 5.2.
The volume form on $`\widehat{}_g`$ induced by the Siegel metric is
(5.5)
$$d\nu _{|\widehat{}_g}=\left(\frac{i}{2}\right)^N\underset{\begin{array}{c}i_N>\mathrm{}>i_1=1\\ j_N>\mathrm{}>j_1=1\end{array}}{\overset{M}{}}\left|\tau _2^1\tau _2^1\right|_{j_1\mathrm{}j_N}^{i_1\mathrm{}i_N}\underset{k=1}{\overset{N}{}}(2\delta _{\mathfrak{1}_{i_k}\mathfrak{2}_{i_k}})\underset{1}{\overset{N}{}}(d\tau _{i_k}d\overline{\tau }_{j_k}),$$
so that
(5.6)
$$\mathrm{Vol}(\widehat{}_g)=_{\widehat{}_g}𝑑\nu _{|\widehat{}_g}.$$
###### Proof.
Let
$$\omega :=\frac{i}{2}\underset{i,j=1}{\overset{M}{}}g_{ij}dZ_id\overline{Z}_j,$$
be the $`(1,1)`$-form associated to the Siegel metric on $`_g`$. By Wirtinger’s theorem , the volume form on a $`d`$-dimensional complex submanifold $`S`$ is
$$\frac{1}{d!}\omega ^d,$$
so that the volume of $`S`$ is expressed as the integral over $`S`$ of a globally defined differential form on $`_g`$. Set $`g_{ij}^\tau :=(2\delta _{\mathfrak{1}_i\mathfrak{2}_i})(\tau _2^1\tau _2^1)_{ij}`$, $`i,jI_M`$, and note that
$`d\nu _{|\widehat{}_g}=`$ $`{\displaystyle \frac{i^N}{2^NN!}}{\displaystyle \underset{\begin{array}{c}i_1,\mathrm{},i_N=1\\ j_1,\mathrm{},j_N=1\end{array}}{\overset{M}{}}}{\displaystyle \underset{k=1}{\overset{N}{}}}g_{i_kj_k}^\tau {\displaystyle \underset{k=1}{\overset{N}{}}}(d\tau _{i_k}d\overline{\tau }_{j_k})`$
$`=`$ $`{\displaystyle \frac{i^N}{2^NN!}}{\displaystyle \underset{\begin{array}{c}i_N<\mathrm{}<i_1=1\\ j_N<\mathrm{}<j_1=1\end{array}}{\overset{M}{}}}{\displaystyle \underset{r,s𝒫_N}{}}ϵ(r)ϵ(s){\displaystyle \underset{k=1}{\overset{N}{}}}g_{i_{r(k)}j_{s(k)}}^\tau {\displaystyle \underset{k=1}{\overset{N}{}}}(d\tau _{i_k}d\overline{\tau }_{j_k}),`$
where $`𝒫_N`$ is the group of permutations of $`N`$ elements and $`ϵ(s)`$ is the sign of the permutation $`s`$. The theorem then follows by the identity
$$\underset{r,s𝒫_N}{}ϵ(r)ϵ(s)\underset{k=1}{\overset{N}{}}g_{i_{r(k)}j_{s(k)}}^\tau =N!\left|\tau _2^1\tau _2^1\right|_{j_1\mathrm{}j_N}^{i_1\mathrm{}i_N}\underset{k=1}{\overset{N}{}}(2\delta _{\mathfrak{1}_{i_k}\mathfrak{2}_{i_k}}).$$
Petri’s basis of $`H^2(K_C^2)`$ corresponds, through the Kodaira-Spencer map, to a basis for the cotangent space of the Teichmüller space. Setting $`d\mathrm{\Xi }_i:=2\pi ik(v_i)`$, $`iI_N`$, it follows by Eq.(4.7) that
(5.7)
$$d\tau _i=\underset{j=1}{\overset{N}{}}w_{ji}d\mathrm{\Xi }_j,iI_M.$$
###### Corollary 5.3.
Fix the points $`p_1,\mathrm{},p_gC`$ in general position, so that $`\{v_i\}_{iI_N}`$, is a basis of $`H^0(K_C^2)`$. The metric on $`\widehat{}_g`$ induced by the Siegel metric is
(5.8)
$$ds_{|\widehat{}_g}^2=\underset{i,j=1}{\overset{N}{}}g_{ij}^\mathrm{\Xi }d\mathrm{\Xi }_id\overline{\mathrm{\Xi }}_j,$$
where $`g_{ij}^\mathrm{\Xi }:=_{k,l=1}^M(2\delta _{\mathfrak{1}_k\mathfrak{2}_k})w_{ik}(\tau _2^1\tau _2^1)_{kl}\overline{w}_{jl}`$.
###### Proof.
Immediate. ∎
By using a suitable basis of $`H^0(K_C^2)`$ and its image under the Kodaira-Spencer map, it turns out that the metric $`g`$ is related to the Bergman reproducing kernel. Fix the points $`z_1,\mathrm{},z_NC`$ satisfying the conditions of Proposition 3.1. The basis $`\{\gamma _i\}_{iI_N}`$ of $`H^0(K_C^2)`$, with $`\gamma _i\gamma _i^2`$, $`iI_N`$, defined by Eq.(3.1) in the case $`n=2`$, satisfies the relations
$$\omega \omega _i=\underset{j=1}{\overset{N}{}}\omega \omega _i(z_j)\gamma _j,v_i=\underset{j=1}{\overset{N}{}}v_i(z_j)\gamma _j,iI_M.$$
Set $`\mathrm{\Gamma }_i:=(2\pi i)^1k(\gamma _i)`$ and $`[v]_{ij}:=v_i(z_j)`$, $`i,jI_N`$.
###### Theorem 5.4.
(5.9)
$$ds_{|\widehat{}_g}^2=\underset{i,j=1}{\overset{N}{}}B^2(z_i,\overline{z}_j)d\mathrm{\Gamma }_id\overline{\mathrm{\Gamma }}_j,$$
in particular, the Siegel induced modular invariant volume form on $`\widehat{}_g`$ is
(5.10)
$$d\nu _{|\widehat{}_g}=\left(\frac{i}{2}\right)^NdetB^2(z_i,\overline{z}_j)_1^N(d\mathrm{\Gamma }_id\overline{\mathrm{\Gamma }}_i),$$
where
$$B(z,\overline{w}):=\underset{i,j=1}{\overset{g}{}}\omega _i(z)(\tau _2^1)_{ij}\overline{\omega }_j(w),$$
$`z,wC`$, is the Bergman reproducing kernel.
###### Proof.
Use $`d\tau _i=_{j=1}^N\omega \omega _i(z_j)d\mathrm{\Gamma }_j`$, $`iI_g`$, and the identity
$$\underset{k,l=1}{\overset{M}{}}(2\delta _{\mathfrak{1}_k\mathfrak{2}_k})\omega \omega _k(z_i)(\tau _2^1\tau _2^1)_{kl}\overline{\omega }\overline{\omega }_l(z_j)=B^2(z_i,\overline{z}_j),i,jI_N.$$
Note that by (5.8) $`_{k,l=1}^N[v]_{ki}g_{kl}^\mathrm{\Xi }[\overline{v}]_{lj}=B^2(z_i,\overline{z}_j)`$, which also follows by (4.7). ∎
## 6. Fay’s trisecant identity from the distinguished basis of $`H^0(K_C^n)`$
Let $`I_i(p):=_{p_0}^p\omega _i`$, $`p_0,pC`$, $`iI_g`$, be the Abel-Jacobi map, which extends to a map from divisors of $`C`$ to the Jacobian $`J(C):=^g/(_g+\tau ^g)`$. We consider Riemann $`\theta `$-functions $`\theta (D+e):=\theta (I(D)+e,\tau )`$, $`eJ(C)`$, evaluated at some $`0`$-degree divisor $`D`$ of $`C`$. By the Riemann vanishing theorem, there is a divisor class $`\mathrm{\Delta }`$ of degree $`g1`$ with $`2\mathrm{\Delta }=K`$, such that $`I(\mathrm{\Delta })`$ is the vector of Riemann constants. Let $`E(z,w)`$ be the prime form, and set
$$\sigma (z):=\mathrm{exp}\left(\underset{i=1}{\overset{g}{}}_{\alpha _i}\omega _i(w)\mathrm{ln}E(z,w)\right).$$
###### Proposition 6.1.
Fix $`n`$ and let $`\{\varphi _i^n\}_{iI_{N_n}}`$ be an arbitrary basis of $`H^0(K_C^n)`$, $`n1`$. Let $`y,x_1,\mathrm{},x_{N_n}`$ be arbitrary points of $`C`$. Then, for $`n=1`$
(6.1)
$$det\varphi _i^1(x_j)=\kappa _1[\varphi ^1]\frac{\theta \left(_1^gx_iy\mathrm{\Delta }\right)_{i<j}^gE(x_i,x_j)_1^g\sigma (x_k)}{\sigma (y)_1^gE(y,x_i)},$$
whereas for $`n>1`$
(6.2)
$$det\varphi _i^n(x_j)=\kappa _n[\varphi ^n]\theta \left(\underset{1}{\overset{N_n}{}}x_i(2n1)\mathrm{\Delta }\right)\underset{i<j}{\overset{N_n}{}}E(x_i,x_j)\underset{i=1}{\overset{N_n}{}}\sigma (x_i)^{2n1},$$
where $`\kappa _1[\varphi ^1]`$ and $`\kappa _n[\varphi ^n]`$ are constants depending only on the choice of the bases.
###### Proof.
$`\kappa _1[\varphi ^1]`$ is a nowhere vanishing section in $`x_j`$, $`jI_g`$, and $`\theta (_1^gx_iy\mathrm{\Delta })=0`$ for $`y=x_1,\mathrm{},x_g`$, so that it is also a nowhere vanishing section in $`y`$ and since it has trivial monodromy it must be a constant. Eq.(6.2) follows by a similar proof. ∎
Set $`w:=_1^{N_n}p_i(2n1)\mathrm{\Delta }`$, $`n>1`$, and assume that $`p_1,\mathrm{},p_{N_n}C`$ satisfy the hypothesis of Proposition 3.1. By Proposition 6.1 we have
(6.3)
$$\gamma _i^n(z)=\frac{\theta (w+zp_i)\sigma (z)^{2n1}_{_{ki}^{k=1}}^{N_n}E(z,p_k)}{\theta (w)\sigma (p_i)^{2n1}_{_{ki}^{k=1}}^{N_n}E(p_i,p_k)},iI_{N_n}.$$
###### Theorem 6.2.
Propositions 3.1 and 6.1 imply the Fay trisecant identity
$$\frac{\theta (w+_{i=1}^m(x_iy_i))_{i<j}E(x_i,x_j)E(y_i,y_j)}{\theta (w)_{i,j}E(x_i,y_j)}=()^{\frac{m(m1)}{2}}det_{ij}\frac{\theta (w+x_iy_j)}{\theta (w)E(x_i,y_j)},$$
$`m2`$, $`x_1,\mathrm{},x_m,y_1,\mathrm{},y_mC`$, $`wJ(C)`$.
###### Proof.
Fix $`m2`$, $`x_1,\mathrm{},x_m,y_1,\mathrm{},y_mC`$ and $`wJ(C)`$, with $`\theta (w)0`$. Choose $`y_1,\mathrm{},y_m`$ distinct, otherwise the identity is trivial. Set $`p_i:=y_i`$, $`iI_m`$, and fix $`n`$, $`N_nm`$, and $`p_{m+1},\mathrm{},p_{N_n}C`$, so that $`w=_1^{N_n}p_i(2n1)\mathrm{\Delta }`$. By Jacobi inversion theorem such a choice is always possible, provided that $`N_nmg`$. It is also clear that, for $`n`$ large enough, $`p_{m+1},\mathrm{},p_{N_n}`$ can be chosen pairwise distinct and distinct from $`y_1,\mathrm{},y_m`$. Eq.(6.2) implies that, by construction, $`det\varphi _i^n(p_j)0`$, for any basis $`\{\varphi _i^n\}_{iI_{N_n}}`$ of $`H^0(K_C^n)`$, since the points $`p_1,\mathrm{},p_{N_n}`$ are pairwise distinct and $`\theta (w)0`$. Therefore, one can define a basis $`\{\gamma _i^n\}_{iI_{N_n}}`$ of $`H^0(K_C^n)`$ by (3.1) and consider $`det\gamma ^n(x_1,\mathrm{},x_m,p_{m+1},\mathrm{},p_{N_n})`$, which can be expressed by (6.3) or by (6.2), with $`\kappa _n[\gamma ^n]`$ determined by applying (6.2) to $`det\gamma _i^n(p_j)=1`$. Comparing these formulas the theorem follows. ∎
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# Transport between edge states in multilayer integer quantum Hall systems: exact treatment of Coulomb interactions and disorder
## I Introduction
Multilayer quantum Hall systems offer a setting in which to study the influence of electron-electron interactions and impurity scattering on tunneling between quantum Hall edge states. Specifically, consider a layered conductor in a magnetic field that is perpendicular to the layers, with the field strength chosen so that a single layer in isolation would have quantised Hall conductance. Then, if interlayer tunneling is not too strong, the multilayer system exhibits a three-dimensional version of the quantum Hall effect and the bulk is insulating at low temperatures. Under these conditions, edge states are present in each layer at the sample surface and are coupled by interlayer tunneling to form a surface phase, which is a chiral, two-dimensional metal.Dohmen ; Balents/Fisher The contribution of this surface phase to the interlayer electron transport properties of such systems has been isolated in experiments on semiconductor multilayers,UCSB1 and is dominant if samples are sufficiently small and cold.
The consequences of impurity scattering for transport in the chiral metal have been discussed extensively from a theoretical viewpoint Dohmen ; Balents/Fisher ; Mathur ; Gruzberg ; Cho ; Plerou ; Sondhi ; Betouras and have been probed experimentally in several ways.UCSB1 ; UCSB0 ; MULTILAYER ; UCSB9 ; Kuraguchi/Osada ; UCSB5 ; UCSB7 ; UCSB2 ; UCSB11 ; UCSB10 ; UCSB12 Crucially, the chiral motion of electrons along the layer edges means that localisation is suppressed.Dohmen ; Balents/Fisher As a result, the surface conductivity in the interlayer direction has a low-temperature limit that is non-zero, even though its measured value may be much smaller than $`e^2/h`$.UCSB1 ; MULTILAYER ; Kuraguchi/Osada Separately, theoretical discussions of conductance fluctationsMathur ; Gruzberg ; Cho ; Plerou ; Betouras have examined both their dependence on geometry in fully phase-coherent samples, and their dependence on the inelastic scattering length when this is smaller than sample size. Observations of reproducible mesoscopic conductance fluctuations,UCSB9 ; UCSB12 induced by small changes of magnetic field within a quantum Hall plateau, demonstrate that interlayer hopping is quantum-mechanically coherent and also provide a way to determine the inelastic scattering length. In addition, magnetoresistance in response to a field component perpendicular to the sample surface has been proposedSondhi and usedUCSB5 ; UCSB7 ; UCSB11 as a method for measuring the elastic scattering length.
In contrast to these studies of disorder effects, past theoretical work on effects due to electron-electron interactions in the chiral metal has been limited. There have been discussions, first, of the temperature dependence of the inelastic scattering lengthBalents/Fisher ; Betouras and, second, of the fact that there is no zero-bias anomaly in the tunneling density of states (or any related contribution to the conductivity), because of ballistic motion of charge in the in-layer direction.Balents/Fisher ; Betouras
Against this background, recent experiments finding a significant temperature dependence to the surface conductivityUCSB2 ; UCSB10 are striking as likely indications of interaction effects, and provide one of the motivations for the work we present here. In particular, the fact that conductivity is observed to increase with increasing temperature presents a puzzle for theory. Some straightforward potential explanations are specifically excluded by the experimental design: large ratios of sample perimeter to cross-sectional area ensure that surface states make the dominant contribution to the measured conductance; and sample perimeters much longer than the inelastic scattering length ensure that weak localisation effects are absent. For samples studied in Ref. UCSB10, , the measured conductivity $`\sigma (T)`$ increases by about $`7\%`$ in the temperature range from $`50`$mK to $`300`$mK, implying a temperature scale of $`\sigma (T)[d\sigma (T)/dT]^14`$K, which is similar to that for other interaction effects in quantum Hall systems
In this paper we study interactions and disorder in the chiral metal, working in the experimentally-relevant limit of weak interlayer tunneling. Treating tunneling perturbatively, Coulomb interactions and impurity scattering can be handled exactly by means of a straighforward application of bosonization. We calculate the full temperature dependence of the conductivity. We also study conductance fluctuations induced by magnetic field changes, obtaining their autocorrelation function and its dependence on temperature. Making appropriate parameter choices, our results for both quantities are consistent with experimental findings. A short account of this work has been presented previously, in Ref. PRL, .
Our work differs from most of the extensive literature on tunneling between quantum Hall edges states in two important ways. First, while much previous work has been concerned with edge states of fractional quantum Hall systems,Wen1 ; Wen2 ; Wen3 ; EdgeReview including multilayer samples,Naud1 ; Naud2 our focus is on the integer quantum Hall effect. Second, whereas most past work (with some exceptions: see Refs. Moon, ; Zulicke, ; Oreg, ; Pryadko, ) has been restricted to systems with only short-range interactions, we find that the long-range nature of Coulomb interactions, which we treat in full, is central for the results we obtain.
The remainder of this paper is organised as follows. We develop a model for the chiral metal in Sec. II and show how bosonization can be used to give an exact description of the collective excitations. Sec. III contains calculations of the temperature dependence of the conductivity. We study conductance fluctuations in Sec. IV, and discuss our results in Sec. V.
## II Modelling the chiral metal
In this section we summarise the physical ingredients that are important for modelling transport between edge states in multilayer conductors and set out the lengthscales that characterise the system. We introduce a Hamiltionian in terms of fermionic operators for edge electrons. We bosonize this Hamiltonian, obtaining a result which is quadratic in boson operators if interlayer tunneling is omitted. Finally, we express the two-electron correlation function that is central to transport calculations in terms of boson correlators.
### II.1 Ingredients, lengthscales, and parameters
A multilayer conductor is illustrated in Fig 1. We use coordinates with the $`x`$-axis parallel to the layer edges, and treat a sample of $`N`$ layers with layer index $`n`$ and layer spacing $`a`$. Consider the system in the presence of a perpendicular magnetic field of strength $`B`$, with the chemical potential lying between the lowest and first excited Landau levels. In the bulk of the sample single particle states at energies close to the chemical potential are localised by disorder. At the sample surface in this energy range, edge states propagate in the confining potential $`V_{\mathrm{edge}}(y)`$ at a velocity $`v`$. Interactions modify the confining potential and the edge velocity: we denote by $`v_\mathrm{F}`$ the velocity allowing for Hartree contributions. Edge states have a width $`w`$ in the $`y`$-direction, which is set by the magnetic length $`l_\mathrm{B}`$ in a clean sample, and by the bulk localisation length $`\xi `$ in the presence of impurities. We use a one-dimensional decription of the edge state in each layer, projected onto the $`x`$-coordinate in the standard way.
Out theoretical treatment takes account only of one edge state in each layer and is therefore appropriate for a system in which electrons are spin polarised. In fact, some of the experiments we refer to, including those on the temperature-dependence of conductivity,UCSB10 are for systems with Landau level filling factor per layer of $`\nu =2`$. It is appropriate to apply our theory to these systems provided electrons with opposite spin directions contribute additively and incoherently to the conductivity.
The system of edge states can be characterised using three lengthscales. First, impurities, which generate only forward scattering with a phase shift, result in an elastic mean free path $`l_{\mathrm{el}}`$, the distance over which a phase shift of order $`2\pi `$ is accumulated. Second, temperature $`T`$ in combination with the velocity $`v_\mathrm{F}`$ can be expressed in terms of the thermal length $`L_\mathrm{T}=\mathrm{}v_\mathrm{F}/k_\mathrm{B}T`$. Third, interlayer tunneling with amplitude $`t_{}`$ can be parameterised by the characteristic distance $`l_{}`$ through which electrons move in the chiral direction between tunneling events. The value of $`l_{}`$ can be expressed in terms of the interlayer diffusion constant $`D`$: since, for small $`t_{}`$, interlayer hops are of length $`a`$ and occur at a rate $`v_\mathrm{F}/l_{}`$, one has $`l_{}=a^2v_\mathrm{F}/D`$. In turn, this can be expressed in terms of the conductivity, using the Einstein relation and the fact that the density of states is $`n=1/2\pi a\mathrm{}v_\mathrm{F}`$, giving $`l_{}=a(e^2/2\pi \mathrm{}\sigma )`$.Betouras
Parameter values for the experiments of Refs. UCSB1, , UCSB11, and UCSB10, are as follows. Samples consist of $`N50`$$`100`$ layers with spacing $`a=30`$nm. The mean free path is estimated UCSB11 to be $`l_{\mathrm{el}}30`$nm. An upper bound on $`v_\mathrm{F}`$, reached in samples with a steep confining potential is $`v_\mathrm{F}\omega _\text{C}l_\mathrm{B}`$, where $`\omega _\text{C}`$ is the cyclotron frequency. It has the value $`\omega _\text{C}l_\text{B}=1.7\times 10^5\text{ms}^1`$ in GaAs at 6.75 T. With this value, $`L_\mathrm{T}10\mu `$m at $`T=100`$mK. Finally, for a surface conductivity of $`\sigma =1.3\times 10^3e^2/2\pi \mathrm{}`$ (which lies within the observed range at $`\nu =2`$), $`l_{}=40\mu `$m. We are therefore concerned with the regime $`l_{\text{el}}L_\text{T}l_{}`$, and this motivates our approach, based on a perturbative treatment of tunneling.
### II.2 Fermionic Hamiltonian
Our model Hamiltonian, $`=_0+_{\text{dis}}+_{\text{hop}}+_{\text{int}}`$, has single-particle terms $`_0`$, $`_{\text{dis}}`$ and $`_{\text{hop}}`$, representing, respectively, free motion along each edge, impurity scattering and interlayer hopping, and a contribution $`_{\text{int}}`$ from Coulomb interactions. We write it in terms of the electron creation operator $`c_{qn}^{}`$ for an edge state with wavevector $`q`$ in layer $`n`$, taking sample perimeter $`L`$ so that $`q=2\pi n_q/L`$, where $`n_q`$ is integer. The creation operator at a point is
$$\psi _n^{}(x)=\frac{1}{\sqrt{L}}\underset{q=\mathrm{}}{\overset{\mathrm{}}{}}e^{iqx}c_{qn}^{}.$$
(1)
We normal order the Hamiltonian with respect to a vacuum in which states are occupied for $`q0`$ and empty otherwise. Then
$`_0`$ $`=`$ $`i\mathrm{}v{\displaystyle \underset{n}{}}{\displaystyle 𝑑x}:\psi _n^{}(x)_x\psi _n(x):,`$ (2)
and
$`_{\text{hop}}`$ $`={\displaystyle \underset{n}{}}{\displaystyle 𝑑x[t_{}\psi _{n+1}^{}(x)\psi _n(x)+\text{H. c.}]}.`$ (3)
The interaction contribution, written in terms of the projected density $`\rho (x)=\psi _n^{}(x)\psi _n(x)`$ with a two-particle potential $`U_{nm}(xx^{})`$, is
$$_{\text{int}}=\frac{1}{2}\underset{nm}{}𝑑x𝑑x^{}:\rho _n(x)U_{nm}(xx^{})\rho _m(x^{}):.$$
(4)
Finally, writing the impurity potential projected onto the edge coordinate in the $`n`$th layer as $`V_n(x)`$, we have
$$_{\text{dis}}=\underset{n}{}𝑑xV_n(x):\psi _n^{}(x)\psi _n(x):.$$
(5)
We take $`V_n(x)`$ to be Gaussian distributed with zero-range correlations and strength $`\mathrm{\Delta }`$: $`[V_n(x)]_{\text{av}}=0`$ and $`[V_n(x)V_n^{}(x^{})]_{\text{av}}=\mathrm{\Delta }\delta _{n,n^{}}\delta (xx^{})`$. This disorder term can be removed by means of a gauge transformation on the fermionic field operators, under which
$$\psi _n^{}(x)e^{i\theta _n(x)}\psi _n^{}(x),$$
(6)
where
$$\theta _n(x)=\frac{1}{\mathrm{}v}_0^x𝑑x^{}V_n(x^{})$$
(7)
is the phase shift acquired under forward scattering from the impurities. The elastic scattering length is related to the disorder strength $`\mathrm{\Delta }`$ by $`l_{\text{el}}=\mathrm{}^2v^2/\mathrm{\Delta }`$. Under this gauge transformation, $`_0+_{\text{dis}}_0`$. The hopping term, however, picks up a dependence on the disorder, and after the transformation is
$$_{\text{hop}}=\underset{n}{}𝑑x[t_{}(n,x)\psi _{n+1}^{}(x)\psi _n(x)+\text{H. c.}],$$
(8)
where
$$t_{}(n,x)=t_{}e^{i(\theta _{n+1}(x)\theta _n(x))}.$$
(9)
We ignore the effects of this gauge transformation on the boundary conditions applying to $`\psi _n(x)`$, which is justified at temperatures large compared to the single-particle level spacing. With this, $`_\text{0}+_{\text{int}}`$ is unaffected by the gauge transformation, and gauge transformed operators $`c_{qn}^{}`$ can be defined by inverting Eq. (1). All further references in this paper to fermionic operators are to the gauge-transformed ones.
### II.3 Bosonised Hamiltonian
We bosonize the Hamiltonian in the standard way, expressing $`_\text{0}+_{\text{int}}`$ in terms of non-interacting collective modes. Since $`_{\text{hop}}`$ transforms into a cosine function of the boson creation and annihilation operators, we treat it perturbatively. To justify this, we require that $`t_{}`$ should be small. Since $`t_{}`$ is a relevant perturbation,Naud1 we also require that temperature should not be too small: $`L_\mathrm{T}l_{}`$.
Boson creation operators are defined in the usual way (see, for example, Ref. vonDelft, ) as
$$b_{qm}^{}=\frac{i}{(n_q)^{1/2}}\underset{r=\mathrm{}}{\overset{\mathrm{}}{}}c_{r+q,m}^{}c_{r,m}^{}$$
(10)
for $`q>0`$. Fourier transforming the interaction potential and expressing the result as a velocity, we introduce
$$u_{nm}(q)=(2\pi \mathrm{})^1𝑑xe^{iqx}U_{nm}(x).$$
(11)
The Fermi velocity renormalised by Hartree interactions is $`v_\mathrm{F}=v_nu_n(0)`$, where the divergence which arises in the sum in the case of Coulomb interactions is cancelled by contributions to $`v`$ from a neutralising background. The Hamiltonian in the absence of hopping (and omitting fermion number terms which appear at electron densities different from that of our vacuum) is
$$_0+_{\text{int}}=\underset{mn}{}\underset{q>0}{}\mathrm{}[v_\mathrm{F}+u_{nm}(q)]qb_{qn}^{}b_{qm}^{}.$$
(12)
The combination $`_0+_{\text{int}}`$ is diagonalised by Fourier transform in the layer index $`n`$. We impose periodic boundary conditions on $`n`$, define the wavevector $`k=2n_k\pi /Na`$, with $`n_k`$ integer and $`\pi /ak<\pi /a`$, and set
$$b_{qk}^{}=\frac{1}{\sqrt{N}}\underset{n=1}{\overset{N}{}}e^{inka}b_{qn}^{},$$
(13)
and
$$u(q,k)=\underset{n}{}e^{inka}u_n(q).$$
(14)
Then
$$_0+_{\text{int}}=\underset{k}{}\underset{q>0}{}\mathrm{}\omega (q,k)b_{qk}^{}b_{qk}^{}$$
(15)
where the excitation frequencies are
$$\omega (q,k)=[v_\text{F}+u(q,k)]q.$$
(16)
The Coulomb interaction, regularised at short distances by a finite width $`w`$ for edge states, has the form
$$U_n(x)=\frac{e^2}{4\pi ϵ_0ϵ_r}\frac{1}{\sqrt{x^2+n^2a^2+w^2}}.$$
(17)
The edge state width $`w`$ is set by the localisation length $`\xi `$ of localised states in the bulk of the sample at the Fermi energy. In a clean sample with well-separated Landau levels, $`\xi l_\mathrm{B}`$, but in a highly disordered sample with Landau levels that are broad in energy one may have $`\xi l_\mathrm{B}`$. The value of $`w`$ proves important in matching our results to experiment, as we discuss in Sec. III.4.
We write the Fourier transform, using the Poisson summation formula, as
$$u(q,k)=v_F\frac{\kappa }{2\pi }\underset{p}{}𝑑x𝑑z\frac{e^{i(qx+kz+2\pi pz/a)}}{\sqrt{x^2+z^2+w^2}}.$$
(18)
and find
$$\omega (q,k)=v_\text{F}q\left(1+\kappa \underset{p}{}Q_p^1e^{wQ_p}\right)$$
(19)
with $`Q_p^2=q^2+(k+2\pi p/a)^2`$ and $`p`$ integer. Here, the inverse screening length $`\kappa e^2/4\pi ϵ_\text{r}ϵ_0\mathrm{}v_\text{F}a`$ characterises the interaction strength.
For isolated layers, taking the limit of large $`a`$, the sum on $`p`$ may be replaced with an integral and one recovers the dispersion relation of edge magnetoplasmons in a single layer system, known from previous work.Volkov1 ; Volkov2
For the multilayer system the expression for the dispersion relation may be simplified in two stages. First, if the layer spacing is small ($`aw`$) the sum on $`p`$ may be omitted, so that
$$\omega (q,k)=v_\text{F}q\left(1+\frac{\kappa e^{w\sqrt{q^2+k^2}}}{\sqrt{q^2+k^2}}\right).$$
(20)
If, in addition, interactions are weak ($`w\kappa ^1`$)
$$\omega (q,k)=v_\text{F}q\left(1+\frac{\kappa }{\sqrt{q^2+k^2}}\right).$$
(21)
In the following we obtain detailed results for systems with wide edges using the dispersion relation of Eq. (20), and for systems with narrow edges using the dispersion relation of Eq. (21).
### II.4 Two-particle correlation function
A central quantity in our calculations of transport properties is the two-fermion correlation function
$$G(x,t)\psi _n^{}(x,t)\psi _{n+1}^{}(x,t)\psi _{n+1}^{}(0,0)\psi _n^{}(0,0),$$
(22)
where $`\mathrm{}\text{Tr}(e^\beta \mathrm{})/\text{Tr}(e^\beta )`$ and operators are written in the Heisenberg representation, with $`𝒪(t)=e^{it/\mathrm{}}𝒪e^{it/\mathrm{}}`$. We evaluate this in the absence of tunneling, so that $`=_0+_{\text{int}}`$.
As a first step, define the boson field operatorfootnote0
$$\varphi _n(x)=\underset{q>0}{}n_q^{1/2}\left(e^{iqx}b_{qn}^{}+e^{iqx}b_{qn}\right)e^{ϵq/2}$$
(23)
where $`ϵ`$ is a short-distance cut-off. Omitting Klein factors (which cancel from $`G(x,t)`$), the fermion and boson field operators are related by
$$\psi _n(x)=(2\pi ϵ)^{1/2}\mathrm{exp}(i\varphi _n(x)).$$
(24)
The correlation function is
$$G(x,t)=\frac{1}{(2\pi ϵ)^2}e^{i\varphi _n(x,t)}e^{i\varphi _{n+1}(x,t)}e^{i\varphi _{n+1}(0,0)}e^{i\varphi _n(0,0)}.$$
(25)
We define its logarithm $`S`$ via
$$G(x,t)\frac{1}{(2\pi )^2}e^S.$$
(26)
Because $``$ is harmonic, $`S`$ can be expressed as
$$\begin{array}{cc}\hfill S=& \frac{1}{2}(\varphi _n(x,t)\varphi _{n+1}(x,t)+\varphi _{n+1}(0,0)\varphi _n(0,0))^2\hfill \\ & +\frac{1}{2}[\varphi _n(x,t)\varphi _{n+1}(x,t),\varphi _n(0,0)\varphi _{n+1}(0,0)]\hfill \\ & 2\mathrm{log}ϵ.\hfill \end{array}$$
(27)
The thermal average and the commutator appearing in this expression can be evaluated in the standard way via a mode expansion, by expressing $`\varphi _n(x,t)`$ in terms of boson creation and annihilation operators using Eq. (23). Taking the thermodynamic limit and so replacing wavevector sums with integrals, with $`\beta =1/k_\mathrm{B}T`$, we arrive at
$`S(x,t,T)=2\mathrm{log}ϵ{\displaystyle \frac{a}{\pi }}{\displaystyle _{\pi /a}^{\pi /a}}𝑑k(1\mathrm{cos}ak){\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dq}{q}}e^{ϵq}`$
$`\times (\mathrm{coth}(\beta \mathrm{}\omega (q,k)/\mathrm{\hspace{0.17em}2})[1\mathrm{cos}(qx\omega (q,k)t)]`$ (28)
$`+i\mathrm{sin}(qx\omega (q,k)t)).`$
It is useful to note that
$$G(x,t)=G(x,t)^{},$$
(29)
and also to define a frequency-dependent correlator,
$$\stackrel{~}{G}(x,\mathrm{\Omega })=𝑑te^{i\mathrm{\Omega }t}G(x,t).$$
(30)
## III Conductivity
In this section we express the conductivity $`\sigma (T)`$ obtained from a Kubo formula in terms of the two-fermion correlation function calculated in Sec. II.4. We also set out the steps required for a numerical evaluation of $`\sigma (T)`$, present our results, and compare them with the experimental data of Ref. UCSB10, .
### III.1 Kubo formula for conductivity
The operator for the interlayer current density between layers $`n`$ and $`n+1`$ is
$$j_n(x)=\frac{ie}{\mathrm{}}\left(t_{}(n,x)\psi _{n+1}^{}(x)\psi _n(x)\text{H. c.}\right).$$
(31)
The real part of the conductivity at frequency $`\mathrm{\Omega }`$ is given by the Kubo formulafootnote
$`\sigma (\mathrm{\Omega },T)={\displaystyle \frac{ia}{\mathrm{}\mathrm{\Omega }L}}{\displaystyle \underset{m}{}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}`$ $`dt\mathrm{sin}\mathrm{\Omega }t{\displaystyle 𝑑x𝑑x^{}}`$
$`\times j_n(x,t)j_m(x^{},0).`$ (32)
To leading order, the interlayer hopping appears only in the current operators, and we evaluate the thermal average using a Hamiltonian from which interlayer hopping is omitted.
Substituting for $`j_n(x,t)`$ using Eq. (31) gives an expression for the conductivity of the chiral metal with a given configuration of disorder: to leading order in $`t_{}(n,x)`$,
$`\sigma (\mathrm{\Omega },T)`$ $`={\displaystyle \frac{2iaL}{\mathrm{}\mathrm{\Omega }}}\left({\displaystyle \frac{e}{\mathrm{}L}}\right)^2{\displaystyle 𝑑x𝑑x^{}_{\mathrm{}}^{\mathrm{}}𝑑t\mathrm{sin}\mathrm{\Omega }t}`$
$`\times t_{}(n,x)t_{}^{}(n,x^{})`$
$`\times \psi _n^{}(x,t)\psi _{n+1}(x,t)\psi _{n+1}^{}(x^{},0)\psi _n(x^{},0).`$ (33)
Averaging over disorder configurations yields
$$[t_{}(n,x)t_{}^{}(n,x^{})]_{\mathrm{av}}=t_{}^2e^{|x|/l_{\mathrm{el}}}$$
(34)
and hence
$`\sigma (\mathrm{\Omega },T)`$ $`={\displaystyle \frac{e^2}{h}}{\displaystyle \frac{8\pi ial_{\text{el}}t_{}^2}{\mathrm{\Omega }\mathrm{}^2}}{\displaystyle \frac{dx}{2l_{\text{el}}}e^{|x|/l_{\text{el}}}_{\mathrm{}}^{\mathrm{}}𝑑t\mathrm{sin}\mathrm{\Omega }t}`$
$`\times \psi _n^{}(x,t)\psi _{n+1}(x,t)\psi _{n+1}^{}(0,0)\psi _n(0,0).`$ (35)
This result can be expressed in terms of the time or frequency dependent two-particle correlation functions defined in Sec. II.4. Setting $`\mathrm{\Omega }=0`$ we find
$`\sigma (T)`$ $`={\displaystyle \frac{e^2}{h}}{\displaystyle \frac{8\pi al_{\text{el}}t_{}^2}{\mathrm{}^2}}{\displaystyle \frac{dx}{2l_{\text{el}}}e^{|x|/l_{\text{el}}}_{\mathrm{}}^{\mathrm{}}𝑑tt\text{Im}G(x,t)}`$ (36)
$`{\displaystyle \frac{e^2}{h}}{\displaystyle \frac{8\pi al_{\text{el}}t_{}^2}{\mathrm{}^2}}{\displaystyle \frac{dx}{2l_{\text{el}}}e^{|x|/l_{\text{el}}}\text{Re}\left[_\mathrm{\Omega }\stackrel{~}{G}(x,\mathrm{\Omega })|_{\mathrm{\Omega }=0}\right]}.`$
For a boson dispersion relation $`\omega (q,k)=v_\mathrm{F}q`$, as results from the Hartree approximation, the fermion correlation function factorises into independent contributions from each layer. These have the form
$$\psi _n^{}(x,t)\psi _n(0,0)=\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}𝑑k\frac{e^{ik(v_\text{F}tx)}}{1+e^{\beta \mathrm{}v_\text{F}k}}$$
(37)
and we find a temperature-independent conductivity
$$\sigma (\mathrm{\Omega },T)=\frac{e^2}{h}\frac{2t_{}^2l_{\text{el}}a}{\mathrm{}^2v_\text{F}^2}\frac{1}{1+\mathrm{\Omega }^2l_{\text{el}}^2/v_\text{F}^2},$$
(38)
which in the zero-frequency limit has the value
$$\sigma _0=\frac{e^2}{h}\frac{2t_{}^2l_{el}a}{\mathrm{}^2v_\text{F}^2}.$$
(39)
More generally, with an arbitrary boson dispersion relation a simplification of Eq. (36) is possible for $`l_{\text{el}}L_\text{T}`$, since $`G(x,t)`$ varies with $`x`$ only on the scale $`L_\text{T}`$ while the correlator $`[t_{}(n,x)t_{}^{}(n,x^{})]_{\mathrm{av}}`$ has range $`l_{\text{el}}`$. We get
$`\sigma (T)`$ $`=4\pi \sigma _0v_\text{F}^2{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑tt\text{Im}G(0,t)`$
$`4\pi \sigma _0v_\text{F}^2\text{Re}\left[_\mathrm{\Omega }\stackrel{~}{G}(0,\mathrm{\Omega })|_{\mathrm{\Omega }=0}\right].`$ (40)
### III.2 Evaluation of $`\sigma (T)`$
To find the temperature dependence of the conductivity we must combine Eqs. (26), (28), and (III.1). A first step before numerical evaluation is to isolate the dependence on the cut-off $`ϵ`$ and take the limit $`ϵ0`$, as we describe in this subsection.
We start from the expression given in Eq. (28) for the logarithm of the two-particle correlation function, which we evaluate at $`x=0`$. It is convenient to separate out a zero-temperature contribution by writing
$$S(t,T)S(t,0)+\mathrm{\Delta }S(t,T)$$
(41)
and also to split $`S(t,0)`$ into real and imaginary parts, with
$$S(t,0)𝒰(t)i𝒱(t),$$
(42)
where $`𝒰(t)`$ and $`𝒱(t)`$ are real for $`t`$ real. Then, writing
$$\sigma (T)=\sigma (0)+\mathrm{\Delta }\sigma (T),$$
(43)
we obtain from Eq. (III.1)
$`\sigma (0)=`$ $`{\displaystyle \frac{2\sigma _0v_\text{F}^2}{\pi }}{\displaystyle _0^{\mathrm{}}}𝑑tte^{𝒰(t)}\mathrm{sin}𝒱(t)`$ (44)
and
$`\mathrm{\Delta }\sigma (T)=`$ $`{\displaystyle \frac{2\sigma _0v_\text{F}^2}{\pi }}{\displaystyle _0^{\mathrm{}}}𝑑tte^{𝒰(t)}\mathrm{sin}𝒱(t)\left[e^{\mathrm{\Delta }S(t,T)}1\right].`$ (45)
In the case of a linear boson dispersion relation, $`\omega (q,k)=v_\mathrm{F}q`$, the functions $`𝒰(t)`$ and $`𝒱(t)`$ have the forms
$`𝒰_{\text{lin}}(t)`$ $`=\mathrm{log}(ϵ^2+v_\text{F}^2t^2)`$ (46)
$`𝒱_{\text{lin}}(t)`$ $`=\pi 2\mathrm{tan}^1(ϵ/v_\text{F}t).`$ (47)
Adding and subtracting these expressions from the ones for $`𝒰(t)`$ and $`𝒱(t)`$ with a general dispersion relation, we find
$`𝒰(t)=𝒰_{\text{lin}}(t)+{\displaystyle \frac{a}{\pi }}{\displaystyle _{\pi /a}^{\pi /a}}𝑑k(1\mathrm{cos}ak)`$
$`\times {\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dq}{q}}e^{ϵq}[\mathrm{cos}(\omega (q,k)t)\mathrm{cos}(v_\text{F}qt)]`$ (48)
and
$`𝒱(t)=𝒱_{\text{lin}}(t)+{\displaystyle \frac{a}{\pi }}{\displaystyle _{\pi /a}^{\pi /a}}𝑑k(1\mathrm{cos}ak)`$
$`\times {\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dq}{q}}e^{ϵq}[\mathrm{sin}(\omega (q,k)t)\mathrm{sin}(v_\text{F}qt)].`$ (49)
Finally, we have
$`\mathrm{\Delta }S(t,T)`$ $`={\displaystyle \frac{a}{\pi }}{\displaystyle _{\pi /a}^{\pi /a}}𝑑k(1\mathrm{cos}ak)`$ (50)
$`\times {\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dq}{q}}`$ $`e^{ϵq}(1\mathrm{cos}\omega (q,k)t)\left[\mathrm{coth}\left({\displaystyle \frac{\beta \mathrm{}\omega (q,k)}{2}}\right)1\right].`$
The advantage of casting the equations for the conductivity in this form is that the momentum integrals in Eqs. (48), (49) and (50) can be performed at $`ϵ=0`$, since the integrands decay fast enough at large $`q`$ for convergence. Dependence on $`ϵ`$ is confined for small $`ϵ`$ to the functions $`𝒰_{\text{lin}}(t)`$ and $`𝒱_{\text{lin}}(t)`$, and from Eqs. (46) and (47) one sees that it is important only for $`t𝒪(ϵ)`$. It is therefore convenient to separate the integration range in Eq. (44) into two parts, $`0t<R`$ and $`Rt<\mathrm{}`$, with $`ϵR1`$. In the first interval $`𝒰(t)=𝒰_{\text{lin}}(t)`$ and $`𝒱=𝒱_{\text{lin}}(t)`$; in the second interval one can set $`ϵ=0`$.
Let the contributions to $`\sigma (0)`$ from the two intervals be $`\sigma ^{(1)}`$ and $`\sigma ^{(2)}`$. Writing $`t^{}=v_\text{F}t/ϵ`$ we have
$$\sigma ^{(1)}=\frac{2ϵ^2\sigma _0}{\pi }_0^{Rϵ^1}𝑑t^{}t^{}e^{𝒰(t^{})}\mathrm{sin}𝒱(t^{})$$
(51)
which gives
$$\sigma ^{(1)}=\frac{2\sigma _0}{\pi }_0^{\mathrm{}}𝑑t^{}t^{}\frac{1}{1+t^2}\frac{2t^{}}{1+t^2}=\sigma _0.$$
(52)
Evaluation of $`\sigma ^{(2)}`$ requires a numerical calculation, and we present results in Sec. III.4.
Finally, turning to the conductivity at non-zero temperature, we note that there are no extra difficulties in the evaluation of $`\mathrm{\Delta }\sigma `$ using Eq. (45). The function $`\mathrm{\Delta }S(t,T)`$, can be computed numerically with $`ϵ=0`$, and $`\mathrm{\Delta }S(t,T)0`$ as $`t0`$, so that $`\mathrm{\Delta }\sigma (T)`$ has no contribution from the integration interval $`0t<R`$ in the limit $`ϵ0`$.
In summary, when evaluating $`\sigma (0)`$ or $`\mathrm{\Delta }\sigma (T)`$ using Eqs. (44) and (45), the functions $`𝒰(t)`$, $`𝒱(t)`$, and $`\mathrm{\Delta }S(t,T)`$ may be evaluated numerically by setting $`ϵ=0`$ in Eqs. (48), (49), and (50), and the results used in Eq. (44) to find $`\sigma ^{(2)}`$. To this one must add $`\sigma ^{(1)}=\sigma _0`$ in order to obtain the zero temperature conductivity $`\sigma (0)`$. These equations combine with Eq. (45) for $`\mathrm{\Delta }\sigma (T)`$ to give a computationally tractable, though non-trivial, expression for $`\sigma (T)`$.
### III.3 Conductivity at zero temperature
The conductivity at zero temperature and zero frequency is determined solely by the low energy limit of the group velocity for excitations, since no other modes are excited as $`T,\mathrm{\Omega }0`$. This zero frequency limit is reached as $`q`$, the wavevector component in the chiral direction, approaches zero. The group velocity, $`\omega (q,k)/q|_{q=0}v_\text{F}\alpha (k)`$, is in general a function of $`k`$, the wavevector component in the interlayer direction.
To determine $`\sigma (0)`$, a useful procedure is to consider a model dispersion relation which is exactly linear in $`q`$: $`\omega (q,k)=v_\text{F}q\alpha (k)`$. A linear dispersion relation is also of interest in its own right. It arises from an interaction that in real space is short range in the chiral direction, $`x`$: $`U_n(x)=g_n\delta (x)`$, giving $`\alpha (k)=1+(2\pi \mathrm{}v_\text{F})^1_ne^{ikna}g_n`$. With a linear dispersion relation, $`q`$-integrals in the expressions leading to $`G(x,t)`$ can be evaluated analytically, greatly simplifying the calculation of conductivity. As we show in the following, for the limit $`l_{\mathrm{el}}L_\mathrm{T}`$ that we consider, a dispersion relation linear in $`q`$ yields a temperature-independent value of conductivity. For interactions, such as the Coulomb potential, that are not short range in $`x`$, linearisation of the dispersion relation gives only an approximation to $`G(x,t)`$. The value of $`\sigma (0)`$ that results from integrating this approximate form for $`G(x,t)`$ is nevertheless exact (at the leading order in $`t_{}`$ considered throughout this paper). This fact is clear on physical grounds, since we have correctly accounted for the dispersion relation at low energy. It may also be derived formally, as follows.
Starting from Eq. (III.1), we deform the contour for the time integral into the semicircle at infinity in the lower half of the complex plane, writing $`t=t_R+\mathrm{i}t_I`$ with $`t_R`$ and $`t_I`$ real. Then in Eq. (28) we have the factor
$$_0^{\mathrm{}}𝑑q\frac{1}{q}\mathrm{exp}(ϵqiqxit_R\omega (q,k)+t_I\omega (q,k)).$$
(53)
This must be evaluated for all values of $`t`$ lying on the deformed time integration contour. When $`|t_R|`$ is large, $`\mathrm{exp}(it_R\omega (q,k))`$ is a rapidly oscillating function of $`q`$, and the $`q`$-integral can be computed using the method of stationary phase: since $`\omega (q,k)`$ is a monotonically increasing function of $`q`$, the dominant contribution comes from the vicinity of the end-point at $`q=0`$. Similarly, when $`t_I`$ is large and negative, $`\mathrm{exp}(t_I\omega (q,k))`$ is small for most values of $`q`$, and the $`q`$-integral can be computed using steepest descents: again, the dominant contribution comes from the vicinity of $`q=0`$. In both instances we may approximate $`\omega (q,k)`$ by its form linearised about $`q=0`$; after linearisation the $`q`$-integral can be evaluated analytically.
This calculation yields
$$\begin{array}{c}\hfill G(0,t)=\frac{1}{(2\pi )^2}\left(\frac{\pi t/\beta \mathrm{}}{\mathrm{sinh}(\pi t/\beta \mathrm{})}\right)^2\frac{1}{v_\text{F}^2}\frac{1}{(ϵ+it)^2}\\ \hfill \times \mathrm{exp}\left(\frac{2a}{\pi }_0^{\pi /a}𝑑k(1\mathrm{cos}ak)\mathrm{log}\alpha (k)\right).\end{array}$$
(54)
Substituting this into Eq. (III.1) we obtain
$`\sigma (T)=`$ $`{\displaystyle \frac{2\sigma _0}{\pi }}\mathrm{exp}\left({\displaystyle \frac{2a}{\pi }}{\displaystyle _0^{\pi /a}}𝑑k(1\mathrm{cos}ak)\mathrm{log}\alpha (k)\right)`$
$`\times {\displaystyle }{\displaystyle \frac{dtϵt^2}{(ϵ^2+t^2)^2}}\left({\displaystyle \frac{\pi t/\beta \mathrm{}}{\mathrm{sinh}(\pi t/\beta \mathrm{})}}\right)^2.`$ (55)
In the limit $`ϵ0`$, the $`t`$ integral gives $`\pi /2`$ regardless of temperature, demonstrating that, for systems with a linear dispersion relation, in the regime $`l_{\mathrm{el}}L_\mathrm{T}`$, $`\sigma (T)`$ is independent of $`T`$. We find
$$\sigma (T)=\sigma _0\mathrm{exp}\left(\frac{2a}{\pi }_0^{\pi /a}𝑑k(1\mathrm{cos}ak)\mathrm{log}\alpha (k)\right).$$
(56)
This is our final result for the dependence of $`\sigma (0)`$ on the dispersion relation as parameterised by $`\alpha (k)`$.
### III.4 Results
We are now in a position to calculate the conductivity for a system with Coulomb interactions by evaluating numerically the formulae we have derived: first, the zero-temperature value using the results from Sec. III.3, and then the full temperature-dependent conductivity using the results from Sec. III.2. We investigate variation of the conductivity with two parameters, the Fermi velocity $`v_\text{F}`$ and the edge state depth $`w`$, and seek values of these parameters for which our results match the experimental data of Ref. UCSB10, . The parameters enter the dispersion relation $`\omega (q,k)`$ directly, and $`v_\text{F}`$ also appears in the inverse screening length $`\kappa `$. The interaction strength is set by the combination $`\kappa a`$ (recall that $`a`$ is the layer spacing). A scale for temperature is set by $`v_\text{F}`$ and $`a`$, via $`T_0\mathrm{}v_\text{F}/ak_\text{B}`$, so that $`T/T_0=a/L_\text{T}`$. A scale for conductivity is given by $`\sigma _0`$, its value in the Hartree approximation.
At a qualitative level, the effect of interactions on the conductivity can be anticipated by starting from the expression given in Eq. (39) for this quantity within the Hartree approximation. In turn, that expression can be understood in terms of a calculation of the interlayer tunneling rate, based on the Fermi golden rule: the rate involves the square of a matrix element between initial and final states on adjacent layers, and a power of the density of states for both the initial and the final states. The squared matrix element, allowing for disorder which affects phases of initial and final states separately, contributes a factor of $`t_{}^2l_{\text{el}}`$ to $`\sigma _0`$. The form of the density of states on a single edge, $`1/2\pi \mathrm{}v_\text{F}`$, implies that $`\sigma _0v_\text{F}^2`$. Returning to a full treatment of the interacting system, we note that the effect of interactions is to generate an energy-dependent group velocity in place of a constant value, $`v_\text{F}`$. In effect, the value of $`\sigma (T)`$ at a particular temperature involves a thermal average of the inverse square of the group velocity. Because Coulomb interactions increase the group velocity at low energy, they decrease conductivity at low temperature; equally, because the group velocity approaches $`v_\text{F}`$ at high energy, the conductivity approaches $`\sigma _0`$ at high temperature.
Turning to detailed results, the dependence of $`\sigma (0)`$ on $`w/a`$ and $`\kappa a`$ is shown in Fig. 2, as obtained from Eq. (56) using $`\alpha (k)=1+\kappa e^{w|k|}/|k|`$. Interactions reduce the value of the conductivity, by a factor which is large if $`\kappa a`$ is large.
The variation of $`\sigma (T)`$ with $`T`$ is illustrated in Fig. 3 for a system with the dispersion relation appropriate for narrow edge states, Eq. (21). In this case the $`k`$ integrals in Eqs. (48) and (49) can be done analytically, leaving only the $`q`$ and $`t`$ integrals to be evaluated numerically.
Finally, the behaviour of $`\sigma (T)`$ for a system with wide edge states ($`wa`$) is presented in Fig. 4. In this case the dispersion relation is as given in Eq. (20), analytical progress does not seem possible, and integrals on $`k`$, $`q`$ and $`t`$ must be evaluated numerically to obtain $`\sigma (T)`$.
We note in passing that we checked that there are only small changes to the results presented when using the more complete form of the interaction given in Eq. (19), including the sum on $`p`$.
Examining these results, it is evident that the general shape of $`\sigma (T)`$ does not vary greatly with parameters: the temperature dependence is quadratic at low temperatures, has a roughly linear region at intermediate temperatures, and approaches $`\sigma _0`$ in the high temperature limit. The quadratic dependence at low temperature is universal, but the extent of the roughly linear region at intermediate temperature is model-dependent. Moreover, scales in this temperature dependence change dramatically with parameter values. The value of the dimensionless temperature $`T/T_0`$ at the crossover between the low and intermediate temperature regimes is dependent on $`\kappa `$ (see Fig. 3) and varies even more strongly with $`w`$ (compare Figs. 3 and 4). In addition, the magnitude of the variation in $`\sigma (T)`$ between low and high $`T`$ depends very much on the values of $`w`$ and $`\kappa a`$. In order to reproduce the experimental observation of a nearly linear increase in $`\sigma (T)`$, by about 7% between the temperatures of 50mK and 300mK,UCSB10 we require parameters which place the experimental temperature window in the intermediate regime for behaviour, so that quadratic variation of $`\sigma (T)`$ with $`T`$ occurs only in a temperature range below 50 mK, and saturation of $`\sigma (T)`$ occurs only above 300mK. Since the available data is not sufficiently detailed to justify a formal fitting procedure, we instead survey the consequences of a range of parameter choices in our results and examine the match to experimental observations.
We begin by considering narrow edges states, using the results shown in Fig. 3. Supposing $`v_\text{F}\omega _\mathrm{C}l_\mathrm{B}`$, which represents an upper bound on $`v_\text{F}`$, we have $`v_\text{F}=1.7\times 10^5\mathrm{ms}^1`$. With $`a=30`$nm, we find $`\kappa a1`$ and $`T_040`$K. Taking these values, the variation in $`\sigma (T)`$ over the experimental temperature range is very small and quadratic, in disagreement with observations. A reduction in the value of $`v_\text{F}`$ serves to decrease the temperature scale $`T_0`$, and also increases $`\kappa `$. It is possible to generate approximately linear variation of $`\sigma (T)`$ with $`T`$ in the experimental temperature range by using a sufficiently small value of $`v_\text{F}`$ (reduced from the upper bound by $`𝒪(10^3)`$), but we know of no reason for $`v_\text{F}`$ to be so small.
We therefore turn to theoretical results for wide edge states, as illustrated in Fig. 4. In this case, we find that large values of $`w`$ greatly reduce the temperature range over which $`\sigma (T)`$ varies quadratically with $`T`$, and can lead to approximately linear variation in the experimental temperature range. A second consequence of large $`w`$ is that the conductivity change $`\sigma (\mathrm{})\sigma (0)`$ is reduced. This tendency can be counteracted by increasing the interaction strength $`\kappa a`$. We find that observed behaviour can be reproduced by taking $`w=4a=120\mathrm{n}\mathrm{m}`$ and $`v_\text{F}=3\times 10^3\text{ms}^1`$ (giving $`\kappa a=50`$). The temperature dependence of $`\sigma (T)`$ obtained using these parameter values is shown in Fig. 5 for temperatures below 400mK.
This choice of parameters, and its implications, merit further discussion. First, we note that there are two separate experimental indications that edge states have a width closer to the value we have adopted, of $`120\mathrm{n}\mathrm{m}`$, than to the conventionally expected value of $`l_\mathrm{B}10\mathrm{n}\mathrm{m}`$. One comes from measurements of bulk hopping transport in multilayer samplesUCSB13 , which give a localisation length of $`\xi =120\mathrm{n}\mathrm{m}`$: one expects $`w\xi `$. The other comes from studies of conductance fluctuations,UCSB12 discussed in Sec. IV. These yield a value for the inelastic scattering length, from the amplitude of fluctuations, and a value for the area of a phase-coherent region perpendicular to the applied field, from the correlation field for fluctuations. The ratio of this phase-coherent area to the inelastic scattering length implies an edge state width which is also much larger than $`l_\mathrm{B}`$: $`w70\mathrm{n}\mathrm{m}`$. Next, turning to the value of $`v_\text{F}`$, which we have taken $`50`$ times smaller than for edge states in a steep confining potential, we note that large edge state width favours a small value for $`v_\text{F}`$, because wide edge states penetrate into the bulk of the sample where both the confining potential gradient and the drift velocity of electrons moving in this potential are small. Finally, we comment on the fact that accepting a small value for $`v_\text{F}`$ implies a large value for $`\sigma _0`$, if other parameters are unchanged. In fact, large $`w`$ acts in the opposite direction, to reduce the effective tunneling amplitude $`t_{}`$ between edge states, since different portions of the edge contribute to the amplitude with different phases, so that there are partial cancellations. To account for the magnitude of the measuredUCSB10 conductivity, $`1.5\times 10^3e^2/2\pi \mathrm{}`$, using the value for the mean free path $`l_{\text{el}}=30\mathrm{n}\mathrm{m}`$ derived from magnetoresistance measurementsUCSB11 requires an effective value of $`t_{}`$ about 50 times smaller than bare estimateUCSB1 of $`0.12`$ meV. This is a surprisingly strong supression of tunneling, though possible if edge states in successive layers have different displacements from the surface, as suggested in Ref. UCSB10, .
## IV Conductance fluctuations
It is found experimentally that mesoscopic fluctuations in the conductance of the chiral metal are induced by small changes of magnetic field within a quantum Hall plateau.UCSB9 ; UCSB12 These conductance fluctuations are observed in samples with a perimeter that is several times larger than the estimated inelastic scattering length. Under such conditions, it is not initially clear why the magnetic field component perpendicular to layers in the sample should influence conductance in this way, since in the simplest picture electron trajectories enclose flux only by encircling the sample. More realistically, a number of possibilities are evident:UCSB12 the sample walls may lie at an angle to the layer normal, either on average or because of surface roughness, or finite edge state width may be important. In our theoretical treatment of conductance fluctuations we avoid specific assumptions about this aspect of the system by considering fluctuations that result from variations in a magnetic field component $`B_{}`$ perpendicular to the sample surface. The amplitude of fluctuations is not affected by this choice. By contrast, the scale for the correlation field of fluctuations is dependent on the model chosen for flux linkage.
In a general setting, there are two possible reasons for the amplitude of conductance fluctuations to decrease with inceasing temperature. One is because of a decrease in the inelastic scattering length; the other is because of thermal smearing. In the case of a chiral metal only the first mechanism operates, because states at different energies are perfectly correlated.Betouras In this sense, conductance fluctuations offer a rather direct probe of interaction effects.
In this section, in place of conductivity $`\sigma `$, we are concerned with the conductance $`g=\sigma L/Na`$ of a finite sample and fluctuations $`\delta g=g[g]_{\mathrm{av}}`$ about its average value. We denote the average within the Hartree approximation by $`g_0\sigma _0L/Na`$. We derive an analytic expression for the autocorrelation function of conductance fluctuations induced by $`B_{}`$. We focus on its temperature dependence at low temperatures, obtaining a scaling form for the regime in which $`\sigma (T)\sigma (0)`$. We compute the scaling function, evaluate our expressions numerically, and compare our results with the observations of Ref. UCSB12, .
### IV.1 Correlation function
The conductance autocorrelation function
$$F(\delta B)=[\delta g(B_{})\delta g(B_{}+\delta B)]_{\text{av}}$$
(57)
is characterised by the amplitude $`F(0)`$ and by the correlation field. An obvious field scale is set by a flux density of one flux quantum $`\mathrm{\Phi }_0`$ through a rectangle with sides proportional to the layer spacing and the thermal length, and we define $`B_0=\mathrm{\Phi }_0/2\pi aL_\text{T}=\mathrm{}/eaL_\mathrm{T}`$. We also introduce a dimensionless field variation $`b=\delta B/B_0`$, which depends on temperature through $`L_\mathrm{T}`$, and a temperature-independent reduced field $`h`$ which has dimensions of wavevector: $`h=b/L_\mathrm{T}e\delta B/a\mathrm{}`$.
With a suitable choice of gauge, the transverse field enters the Hamiltonian only as a phase for interlayer hopping. Taking for convenience $`B_{}=0`$, in the presence of non-zero $`\delta B`$ Eq. (9) is modified to
$`t_{}(n,x)=t_{}e^{i(\theta _{n+1}(x)\theta _n(x)+hx)}.`$ (58)
This additional, field-dependent phase alters $`_{\text{hop}}`$ and consequently the current operator.
An expression for the conductance of a sample with a specific disorder configuration is obtained by scaling Eq. (III.1) with the sample dimensions. Taking account of the field-dependent phases in the current operator and substituting into the definition of $`F(\delta B)`$, after some manipulation we arrive at
$`F(\delta B)`$ $`={\displaystyle \frac{g_0^2\pi ^2v_\text{F}^4}{L^2l_{\text{el}}^2N^2}}{\displaystyle \underset{n,m}{}}{\displaystyle 𝑑x𝑑x^{}𝑑y𝑑y^{}}`$ (59)
$`\times `$ $`{\displaystyle 𝑑titG(xx^{},t)𝑑t^{}it^{}G(yy^{},t^{})}`$
$`\times `$ $`\left(e^{ih(xx^{})}+e^{ih(xx^{})}\right)e^{C(x,x^{})}e^{C(y,y^{})}`$
$`\times `$ $`\left(e^{D_{nm}(x,x^{};y,y^{})}+e^{D_{nm}(x,x^{};y,y^{})}2\right).`$
Two contributions to this expression arise from the disorder average:
$$C(x,x^{})=\frac{1}{2}[(\theta _{n+1}(x)\theta _n(x)\theta _{n+1}(x^{})+\theta _n(x^{}))^2]_{\text{av}}$$
(60)
and
$`D_{nm}`$ $`(x,x^{};y,y^{})=[(\theta _{n+1}(x)\theta _n(x)\theta _{n+1}(x^{})+\theta _n(x^{}))`$
$`\times `$ $`(\theta _{m+1}(y)\theta _m(y)\theta _{m+1}(y^{})+\theta _m(y^{}))]_{\text{av}}.`$ (61)
Both may be evaluated using the result (for $`x,y>0`$)
$$[\theta _n(x)\theta _m(y)]_{\text{av}}=\frac{\delta _{nm}}{l_{\text{el}}}\mathrm{min}\{x,y\}.$$
(62)
The equation for $`C`$ gives
$$e^{C(x,x^{})}=e^{|xx^{}|/l_{\text{el}}},$$
(63)
which in the limit of small $`l_{\text{el}}`$ can be written $`2l_{\text{el}}\delta (xx^{})`$. The expression for $`D`$ is more complicated: one finds
$`D_{nm}(x,x^{};y,y^{})={\displaystyle \frac{R(x,x^{};y,y^{})}{l_{\text{el}}}}\left(2\delta _{nm}\delta _{n+1,m}\delta _{n1,m}\right).`$ (64)
The function $`R(x,x^{};y,y^{})`$ gives the overlap between the two directed intervals on the real line $`xx^{}`$ and $`yy^{}`$: for example, $`R(1,5;4,9)=R(5,1;4,9)=1`$. On substituting these expressions for $`C`$ and $`D`$ into Eq.(59), we obtain
$`F(b)={\displaystyle \frac{g_0^2\pi ^2v_\text{F}^4}{L^2l_{\text{el}}^2N}}{\displaystyle 𝑑x𝑑x^{}𝑑y𝑑y^{}\left(e^{ih(xx^{})}+e^{ih(xx^{})}\right)}`$
$`\times {\displaystyle }dtitG(xx^{},t){\displaystyle }dt^{}it^{}G(yy^{},t^{})e^{|xx^{}|/l_{\text{el}}}`$
$`\times e^{|yy^{}|/l_{\text{el}}}\{e^{2R(x,x^{};y,y^{})/l_{\text{el}}}+e^{2R(x,x^{};y,y^{})/l_{\text{el}}}2`$
$`+\mathrm{\hspace{0.17em}2}e^{R(x,x^{};y,y^{})/l_{\text{el}}}+2e^{R(x,x^{};y,y^{})/l_{\text{el}}}4\}.`$ (65)
Examining where the weight of the integrand lies with respect to the spatial integrals in Eq. (65), one sees that the term in braces vanishes except in places where $`R0`$. We consider different types of contributions from these regions, and keep only those which are leading order for $`L_\text{T}l_{\text{el}}`$. First, consider regions in which $`|xy|l_{\text{el}}`$ but $`|xx^{}|l_{\text{el}}`$. The small factor $`e^{|xx^{}|/l_{\text{el}}}`$ is compensated by the first term in the braces if $`|x^{}y^{}|l_{\text{el}}`$. Then
$`e^{|xx^{}|/l_{\text{el}}}e^{|yy^{}|/l_{\text{el}}}e`$ $`{}_{}{}^{2R(x,x^{};y,y^{})/l_{\text{el}}}=`$ (66)
$`e^{(|xy||x^{}y^{}|)/l_{\text{el}}}.`$
Since $`G(x,t)`$ has a range in $`x`$ of order $`L_\text{T}`$, the resulting contribution to $`F(\delta B)`$ is $`𝒪(L_\text{T}/L)`$. Another contribution of the same order arises from regions where $`|xy^{}|l_{\text{el}}`$ and $`|x^{}y|l_{\text{el}}`$. Subleading contributions come from regions where all four spatial variables are within an elastic length of one another. These contributions are $`𝒪(l_{\text{el}}/L)`$.
Keeping only the leading order terms, the expression for the correlation function has the much simplified form
$`F(\delta B)`$ $`={\displaystyle \frac{4g_0^2\pi ^2v_\text{F}^4}{NL}}{\displaystyle 𝑑x(e^{ihx}+e^{ihx})𝑑tit𝑑t^{}it^{}}`$
$`\times \left(G(x,t)G(x,t^{})+G(x,t)G(x,t^{})\right).`$ (67)
Using the symmetry of $`G(x,t)`$ (see Eq. (29)) one finds
$$F(\delta B)=\frac{g_0^2}{NL}_{\mathrm{}}^{\mathrm{}}𝑑xe^{ihx}[f(x)]^2,$$
(68)
where
$$f(x)4\pi v_\text{F}^2_{\mathrm{}}^{\mathrm{}}𝑑tt\text{Im}G(x,t).$$
(69)
### IV.2 Computing the correlation function
In order to compare our theory for conductance fluctuations with experiment, we need to be able to calculate $`F(\delta B)`$ for various values of the temperature and parameters $`v_\text{F}`$ and $`w`$. Although it is possible to use a computer to evaluate the form of $`F(B_{})`$ given in Eq. (68) without further approximation, it is far easier to make progress by calculating $`G(x,t)`$ for a linearised dispersion relation. This approach is exact in the low-temperature regime defined by the condition $`\sigma (T)\sigma (0)`$, and we proceed to use it in our calculations.
In the low temperature regime where the linearised dispersion relation may be used, $`F(B_{})`$ has a scaling form. To make this apparent, it is helpful to recast equations in terms of dimensionless variables, characterising $`\delta B`$ by $`b`$ in place of $`h`$, and introducing $`\widehat{x}=x/L_T`$ and $`\widehat{t}=v_\text{F}t/L_\text{T}`$. Writing $`G(x,t)=(2\pi L_\text{T})^2\widehat{G}(\widehat{x},\widehat{t})`$ and $`f(L_\text{T}\widehat{x})=\widehat{f}(\widehat{x})`$, for a linear dispersion relation, $`\omega (q,k)=qv_\text{F}\alpha (k)`$, we have
$`\widehat{G}(\widehat{x},\widehat{t})=\mathrm{exp}\{{\displaystyle \frac{2a}{\pi }}{\displaystyle _0^{\pi /a}}dk(1\mathrm{cos}ak)`$
$`\times [\mathrm{log}|\widehat{x}\alpha (k)\widehat{t}|\mathrm{log}\left({\displaystyle \frac{\pi [\alpha (k)\widehat{t}\widehat{x}]/\alpha (k)}{\mathrm{sinh}(\pi [\alpha (k)\widehat{t}\widehat{x}]/\alpha (k))}}\right)]\}`$
$`\times \mathrm{exp}\left\{ia{\displaystyle _0^{\pi /a}}𝑑k(1\mathrm{cos}ak)\text{sgn}(\widehat{x}\alpha (k)\widehat{t})\right\}`$ (70)
and
$$\widehat{f}(\widehat{x})=\frac{1}{\pi }_{\mathrm{}}^{\mathrm{}}𝑑\widehat{t}\widehat{t}\text{Im}\{\widehat{G}(\widehat{x},\widehat{t})\}.$$
(71)
Then the conductance autocorrelation function has the form
$$F(\delta B)=\frac{g_0^2L_\text{T}}{NL}C\left(\delta B/B_0\right)$$
(72)
with scaling function
$$C(b)=_{\mathrm{}}^{\mathrm{}}𝑑\widehat{x}e^{ib\widehat{x}}[\widehat{f}(\widehat{x})]^2.$$
(73)
In this form $`F(\delta B)`$ depends on temperature $`T`$ and magnetic field difference $`\delta B`$ only through the scaling variables $`L_\text{T}/L`$ and $`\delta B/B_0`$. The thermal length $`L_\text{T}`$ plays the role of an inelastic scattering length, in the sense that it determines both the amplitude of conductance fluctuations and (through $`B_0`$) their correlation field. Such behaviour is initally surprising, since $`L_\text{T}`$ is independent of interaction strength. In fact, of course, the form of the scaling function $`C(b)`$ depends parametrically on interaction strength.
For weak interactions this dependence of $`C(b)`$ on $`\kappa `$ can be extracted analytically, as follows. First, note from Eq. (21) that $`\alpha (k)=1+\kappa /|k|`$. Also, in Eqs. (70), (71) and (73), change variables from $`\widehat{x},\widehat{t}`$ to $`y,p`$ with $`\widehat{x}=y/\kappa `$ and $`\widehat{t}=yp+y/\kappa `$. Then
$`\underset{\kappa 0}{lim}`$ $`\widehat{G}(y/\kappa ,p+y/\kappa )g(y,p)`$
$`=`$ $`\mathrm{exp}\{{\displaystyle \frac{2a}{\pi }}{\displaystyle _0^{\pi /a}}dk(1\mathrm{cos}ak)`$
$`\times [\mathrm{log}|y(p+1/k)|\mathrm{log}\left({\displaystyle \frac{\pi y[p+1/k]}{\mathrm{sinh}(\pi y[p+1/k])}}\right)]\}`$
$`\times \mathrm{exp}\left\{ia{\displaystyle _0^{\pi /a}}𝑑k(1\mathrm{cos}ak)\text{sgn}(y[p+1/k])\right\}`$
and
$$\underset{\kappa 0}{lim}\widehat{f}(y/\kappa )\stackrel{~}{f}(y)=\frac{y^2}{\pi }_{\mathrm{}}^{\mathrm{}}𝑑pp\text{Im}\{g(y,p)\}.$$
The $`\kappa `$-dependence of the scaling function is hence isolated for small $`\kappa `$ as
$$C(b)=\frac{1}{\kappa }_{\mathrm{}}^{\mathrm{}}𝑑y\mathrm{exp}(iyb/\kappa )[\stackrel{~}{f}(y)]^2,$$
(74)
demonstrating that the amplitude of conductance fluctuations grows and that the correlation field shrinks as interactions are made weaker. In both cases, the variation implies an inelastic scattering length that diverges as $`\kappa ^1`$ for weak interactions. Such a dependence of the inelastic scattering length on interaction strength is long-established in non-chiral, one-dimensional conductors.Apel
In order to find the form of the scaling function and to study its $`\kappa `$-dependence at general $`\kappa `$, a three-dimensional numerical integration is necessary. We compute $`\widehat{G}(\widehat{x},\widehat{t})`$, then $`\widehat{f}(\widehat{x})`$, and then the scaling function $`C(b)`$ itself.
### IV.3 Results
We illustrate the form of the scaling function $`C(\delta B/B_0)`$ for a range of parameter values in a sequence of three figures. Its dependence on interaction strength $`\kappa a`$ is shown for narrow edge states in Fig. 6 and for $`w=a`$ in Fig. 7. In both cases, smaller interaction strength leads to a larger amplitude for conductance fluctuations and a smaller correlation field, as may be anticipated on the grounds that weaker interactions lead to a longer inelastic scattering length. In Fig. 8 $`C(\delta B/B_0)`$ is shown for $`\kappa =50`$ and $`w=4a`$, the parameter values suggested by the comparison of our conductivity calculations with experiment. We discuss experimental data on conductance fluctuations in Sec. IV.4.
Finally, the increase in the amplitude of conductance fluctuations with dereasing $`\kappa `$ is illustrated in Fig. 9.
### IV.4 Comparison with experiment and previous theory
The exact treatment of disorder and interactions provided by the calculations we have decribed presents an opportunity to test the standard theoretical treatment of conductance fluctuations, in which a single inelastic scattering length $`l_{\mathrm{in}}`$, or equivalently a scattering rate $`v_\text{F}/l_{\mathrm{in}}`$ is used as a cut-off in perturbation theory. For the chiral metal, such calculations have been described in Ref. Betouras, . They yield a Lorentzian scaling function
$$F(\delta B)=\frac{2g_0^2}{NL}\frac{l_{\text{in}}}{1+z^2}$$
(75)
with $`z=2\pi \delta Bl_{\text{in}}a/\mathrm{\Phi }_0`$. A comparison between the functional form we obtain for $`F(\delta B)`$ and a Lorenztian is given in Fig. 8: while the two functions are similar, the discrepancies are worth attention because they indicate behaviour which cannot be characterised by a single relaxation time. A similar comparison can be made in the Fourier transformed domain, in terms of the function $`f(x)`$. To reproduce Eq. (75) from our Eq. (68), we would require $`l_{\text{in}}=L_\text{T}`$ and
$$\widehat{f}(\widehat{x})=e^{|\widehat{x}|/2},$$
(76)
where exponential decay is indicative of a single lifetime $`l_{\text{in}}/v_\text{F}`$ for excitations. The form we obtain for $`\widehat{f}(\widehat{x})`$ is shown in Fig. 10. The absence of a cusp at $`x=0`$ indicates that there is of a range of relaxation times in the system. In addition, the fact that $`f(0)1`$ is an interaction effect (from Eq. (III.1) one sees that $`f(0)=\sigma (0)/\sigma _0`$) not allowed for in the standard perturbative treatment.
We close this section with a comparison between the experiments of Ref. UCSB12, and our results, using the same parameters, $`\kappa a=50`$ and $`w=4a`$, that provided a match for the behaviour of $`\sigma (T)`$. For the experimental base temperature of $`T=70\text{mK}`$, we use our approach to determine the amplitude of conductance fluctuations. As a way to present the result, we then follow the experimental analysisUCSB12 in using Eq. (75) to obtain a value for $`l_{\text{in}}`$ of $`0.3\mu \text{m}`$. The experimental value, extracted in the same way, is $`l_{\text{in}}1\mu \text{m}`$. Since the calculated amplitude of conductance fluctuations varies by several orders of magnitude over the range of parameter values we have investigated, and since no new adjustment of parameters was involved in our discussion of conductance fluctuations, we find the rough agreement between these two values of $`l_{\text{in}}`$ very encouraging.
## V Conclusions
In summary, for the system of weakly coupled quantum Hall edge states that we have studied, bosonisation provides a very complete treatment of the interplay between electron-electron interactions and disorder. We have shown that interaction effects can account for the observed temperature dependence of interlayer conductivity, provided we allow for finite edge state width and adopt a value for the edge state velocity that is rather smaller than previously supposed. We have investigated conductance fluctuations within the same theoretical approach, showing how they are suppressed with increasing temperature, with a characteristic lengthscale $`L_\text{T}T^1`$. Encouragingly, the same parameter values used to match the measured behaviour of conductivity reproduce approximately the observed fluctuation amplitude. From a theoretical viewpoint, it is interesting that such dephasing effects can be generated from a description based on harmonic collective modes, simply via the nonlinear relation between boson and fermion operators.
###### Acknowledgements.
We thank E. G. Gwinn for very helpful discussions and J. J. Betouras for previous collaborations. The work was supported by EPSRC under Grant GR/R83712/01 and by the Dutch FOM foundation.
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# 1 Introduction
## 1 Introduction
The long gluino lifetime is a trademark of Split Supersymmetry . The experimental discovery of a slowly–decaying gluino would not only be a strong indication for Split Supersymmetry, but it would also allow for a measurement of the effective supersymmetry–breaking scale $`\stackrel{~}{m}`$, which cannot be directly extracted from particle dynamics at the LHC. Moreover, the gluino lifetime is a crucial parameter to determine the cosmological constraints on the theory . Therefore, for both experimental and theoretical considerations, it is very important to have a precise prediction of the gluino lifetime and branching ratios.
For what concerns the gluino decay processes in the MSSM, tree–level results for the decays into chargino or neutralino and two quarks and one–loop results for the radiative decay into neutralino and gluon can be found in the literature . In Split Supersymmetry, however, the quantum corrections to the gluino decay processes can be very significant, because they are enhanced by the potentially large logarithm of the ratio between the gluino mass $`m_{\stackrel{~}{g}}`$ and the scale $`\stackrel{~}{m}`$ at which the interactions responsible for gluino decay are mediated. A fixed–order calculation of these processes in Split Supersymmetry would miss terms that are enhanced by higher powers of the large logarithm. In order to get a reliable prediction for the gluino decay width, the large logarithmic corrections have to be resummed by means of standard renormalization group techniques.
Recently, a calculation of the gluino decay widths in Split Supersymmetry was presented in ref. , working at tree level for 3–body decays and in (not resummed) one–loop approximation for 2–body decays. In this paper we will present a calculation of the gluino decay processes that includes all leading corrections in $`\alpha _s`$ and $`\alpha _t`$, the strong and top–Yukawa coupling constants. As we will show, the inclusion and resummation of leading–order corrections give sizeable modifications of the gluino branching ratios, even for moderate values of $`\stackrel{~}{m}`$.
The structure of the paper is as follows: in sect. 2 we list the operators in the effective Lagrangian of Split Supersymmetry that are responsible for the decays of the gluino, and the high–energy boundary conditions on the corresponding Wilson coefficients; in sect. 3 we determine the renormalization group evolution of the Wilson coefficients, and we express the operators in the low–energy effective Lagrangian in terms of mass eigenstates; in sect. 4 we discuss our numerical results for the branching ratios and total width of the gluino decays in Split Supersymmetry; in sect. 5 we consider the possibility of gluinos decaying into gravitino; in sect. 6 we present our conclusions. Finally, in the appendix we provide the analytical formulae for the gluino decay widths.
## 2 The Effective Lagrangian
Below the squark and slepton mass scale $`\stackrel{~}{m}`$, the effective Lagrangian of Split Supersymmetry describes the dynamics of Standard Model (SM) particles together with higgsinos and gauginos. At the level of renormalizable interactions, there is a conserved $`G`$–parity (under which only the gluino is odd) preventing gluino decay. However, integrating the squarks out of the underlying supersymmetric theory induces non–renormalizable interactions that violate the $`G`$–parity. Restricting our analysis up to dimension–6 operators, the $`G`$–odd effective Lagrangian at the matching scale $`\stackrel{~}{m}`$ is given by
$$=\frac{1}{\stackrel{~}{m}^2}\underset{i=1}{\overset{7}{}}C_i^{\stackrel{~}{B}}Q_i^{\stackrel{~}{B}}+\frac{1}{\stackrel{~}{m}^2}\underset{i=1}{\overset{2}{}}C_i^{\stackrel{~}{W}}Q_i^{\stackrel{~}{W}}+\frac{1}{\stackrel{~}{m}^2}(\underset{i=1}{\overset{5}{}}C_i^{\stackrel{~}{H}}Q_i^{\stackrel{~}{H}}+\mathrm{h}.\mathrm{c}.).$$
(1)
We are working in the basis of interaction eigenstates for gauginos and higgsinos, neglecting the effect of electroweak symmetry breaking, since $`\stackrel{~}{m}M_Z`$. The $`G`$–odd operators involving the $`B`$–ino ($`\stackrel{~}{B}`$) are
$`Q_1^{\stackrel{~}{B}}`$ $`=`$ $`\overline{\stackrel{~}{B}}\gamma ^\mu \gamma _5\stackrel{~}{g}^a{\displaystyle \underset{k=1}{\overset{2}{}}}\overline{q}_L^{(k)}\gamma _\mu T^aq_L^{(k)}`$ (2)
$`Q_2^{\stackrel{~}{B}}`$ $`=`$ $`\overline{\stackrel{~}{B}}\gamma ^\mu \gamma _5\stackrel{~}{g}^a{\displaystyle \underset{k=1}{\overset{2}{}}}\overline{u}_R^{(k)}\gamma _\mu T^au_R^{(k)}`$ (3)
$`Q_3^{\stackrel{~}{B}}`$ $`=`$ $`\overline{\stackrel{~}{B}}\gamma ^\mu \gamma _5\stackrel{~}{g}^a{\displaystyle \underset{k=1}{\overset{2}{}}}\overline{d}_R^{(k)}\gamma _\mu T^ad_R^{(k)}`$ (4)
$`Q_4^{\stackrel{~}{B}}`$ $`=`$ $`\overline{\stackrel{~}{B}}\gamma ^\mu \gamma _5\stackrel{~}{g}^a\overline{q}_L^{(3)}\gamma _\mu T^aq_L^{(3)}`$ (5)
$`Q_5^{\stackrel{~}{B}}`$ $`=`$ $`\overline{\stackrel{~}{B}}\gamma ^\mu \gamma _5\stackrel{~}{g}^a\overline{t}_R\gamma _\mu T^at_R`$ (6)
$`Q_6^{\stackrel{~}{B}}`$ $`=`$ $`\overline{\stackrel{~}{B}}\gamma ^\mu \gamma _5\stackrel{~}{g}^a\overline{b}_R\gamma _\mu T^ab_R`$ (7)
$`Q_7^{\stackrel{~}{B}}`$ $`=`$ $`\overline{\stackrel{~}{B}}\sigma ^{\mu \nu }\gamma _5\stackrel{~}{g}^aG_{\mu \nu }^a,`$ (8)
where $`k`$ is a generation index, $`T^a`$ are the SU(3) generators and $`G_{\mu \nu }^a`$ is the gluon field strength. Assuming that the squark mass matrices are flavour–diagonal, the Wilson coefficients of the operators $`Q_i^{\stackrel{~}{B}}`$ at the matching scale $`\stackrel{~}{m}`$ are
$`C_1^{\stackrel{~}{B}}(\stackrel{~}{m})=C_4^{\stackrel{~}{B}}(\stackrel{~}{m})={\displaystyle \frac{g_sg^{}}{6}}r_{\stackrel{~}{q}_L},`$ $`C_2^{\stackrel{~}{B}}(\stackrel{~}{m})=C_5^{\stackrel{~}{B}}(\stackrel{~}{m})={\displaystyle \frac{2g_sg^{}}{3}}r_{\stackrel{~}{u}_R},`$ (9)
$`C_3^{\stackrel{~}{B}}(\stackrel{~}{m})=C_6^{\stackrel{~}{B}}(\stackrel{~}{m})={\displaystyle \frac{g_sg^{}}{3}}r_{\stackrel{~}{d}_R},`$ $`C_7^{\stackrel{~}{B}}(\stackrel{~}{m})={\displaystyle \frac{g_s^2g^{}}{128\pi ^2}}(m_{\stackrel{~}{g}}m_{\stackrel{~}{B}}){\displaystyle \underset{q}{}}(r_{\stackrel{~}{q}_L}r_{\stackrel{~}{q}_R})Q_q,`$ (10)
where $`r_{\stackrel{~}{q}}=\stackrel{~}{m}^2/m_{\stackrel{~}{q}}^2`$. Note that $`C_7^{\stackrel{~}{B}}`$ vanishes for mass–degenerate squarks.
The $`G`$–odd operators involving the $`W`$–ino ($`\stackrel{~}{W}`$) are
$`Q_1^{\stackrel{~}{W}}`$ $`=`$ $`\overline{\stackrel{~}{W}^A}\gamma ^\mu \gamma _5\stackrel{~}{g}^a{\displaystyle \underset{k=1}{\overset{2}{}}}\overline{q}_L^{(k)}\gamma ^\mu \tau ^AT^aq_L^{(k)}`$ (11)
$`Q_2^{\stackrel{~}{W}}`$ $`=`$ $`\overline{\stackrel{~}{W}^A}\gamma ^\mu \gamma _5\stackrel{~}{g}^a\overline{q}_L^{(3)}\gamma ^\mu \tau ^AT^aq_L^{(3)},`$ (12)
where $`\tau ^A`$ are the Pauli matrices. The matching conditions for the Wilson coefficients are
$$C_1^{\stackrel{~}{W}}(\stackrel{~}{m})=C_2^{\stackrel{~}{W}}(\stackrel{~}{m})=\frac{g_sg}{2}r_{\stackrel{~}{q}_L}.$$
(13)
For the higgsinos, we use a compact notation in which the two Weyl states $`\stackrel{~}{H}_u`$ and $`\stackrel{~}{H}_d`$ are combined in a single Dirac fermion $`\stackrel{~}{H}\stackrel{~}{H}_u+\epsilon \stackrel{~}{H}_d^c`$, where $`\epsilon `$ is the antisymmetric matrix (with $`\epsilon _{12}=1`$) acting on the SU(2) indices. The states $`\stackrel{~}{H}_u`$ and $`\stackrel{~}{H}_d`$ can be recovered by chiral decomposition, $`\stackrel{~}{H}_u=\stackrel{~}{H}_L`$ and $`\stackrel{~}{H}_d=\epsilon (\stackrel{~}{H}^c)_L`$. Keeping only the third–generation Yukawa couplings, the $`G`$–odd operators involving higgsinos are
$`Q_1^{\stackrel{~}{H}}`$ $`=`$ $`\overline{\stackrel{~}{H}}_L\stackrel{~}{g}_R^a\epsilon \overline{q}_L^{(3)}T^at_R`$ (14)
$`Q_2^{\stackrel{~}{H}}`$ $`=`$ $`\overline{\stackrel{~}{H}}_L\sigma ^{\mu \nu }\stackrel{~}{g}_R^a\epsilon \overline{q}_L^{(3)}\sigma _{\mu \nu }T^at_R`$ (15)
$`Q_3^{\stackrel{~}{H}}`$ $`=`$ $`\overline{\stackrel{~}{H}}_R\stackrel{~}{g}_L^a\overline{b}_RT^aq_L^{(3)}`$ (16)
$`Q_4^{\stackrel{~}{H}}`$ $`=`$ $`\overline{\stackrel{~}{H}}_R\sigma ^{\mu \nu }\stackrel{~}{g}_L^a\overline{b}_R\sigma _{\mu \nu }T^aq_L^{(3)}`$ (17)
$`Q_5^{\stackrel{~}{H}}`$ $`=`$ $`\overline{\stackrel{~}{H}}_L\sigma ^{\mu \nu }\stackrel{~}{g}_R^ahG_{\mu \nu }^a,`$ (18)
where $`h`$ is the Higgs doublet. The Wilson coefficients at the matching scale $`\stackrel{~}{m}`$ are
$$C_1^{\stackrel{~}{H}}(\stackrel{~}{m})=\frac{g_sh_t}{\sqrt{2}\mathrm{sin}\beta }(r_{\stackrel{~}{q}_L}r_{\stackrel{~}{u}_R}),C_2^{\stackrel{~}{H}}(\stackrel{~}{m})=\frac{g_sh_t}{4\sqrt{2}\mathrm{sin}\beta }(r_{\stackrel{~}{q}_L}+r_{\stackrel{~}{u}_R}),$$
(19)
$$C_3^{\stackrel{~}{H}}(\stackrel{~}{m})=\frac{g_sh_b}{\sqrt{2}\mathrm{cos}\beta }(r_{\stackrel{~}{q}_L}r_{\stackrel{~}{d}_R}),C_4^{\stackrel{~}{H}}(\stackrel{~}{m})=\frac{g_sh_b}{4\sqrt{2}\mathrm{cos}\beta }(r_{\stackrel{~}{q}_L}+r_{\stackrel{~}{d}_R}),$$
(20)
$$C_5^{\stackrel{~}{H}}(\stackrel{~}{m})=\frac{g_s^2h_t^2}{32\sqrt{2}\pi ^2\mathrm{sin}\beta }(r_{\stackrel{~}{q}_L}+r_{\stackrel{~}{u}_R}).$$
(21)
Here $`h_t`$ and $`h_b`$ are the top and bottom Yukawa couplings, and $`\mathrm{tan}\beta `$ is a free parameter of Split Supersymmetry.
Before proceeding to the operator renormalization, we want to make some remarks.
(i) We recall that all coupling constants appearing in the expressions of the Wilson coefficients given above have to be computed at the scale $`\stackrel{~}{m}`$.
(ii) Note that we have given the Wilson coefficients of the 4–fermion operators at the leading perturbative order, while the coefficients of the operators $`Q_7^{\stackrel{~}{B}}`$ and $`Q_5^{\stackrel{~}{H}}`$ are given at the next order (one–loop approximation). The operator anomalous dimensions will be computed in sect. 3 at the leading order in the strong and top–Yukawa couplings $`\alpha _s=g_s^2/(4\pi )`$ and $`\alpha _t=h_t^2/(4\pi )`$. Therefore, the gluino 3–body decays, mediated only by 4–fermion operators, will be computed by resumming all $`\alpha _{s,t}\mathrm{ln}(\stackrel{~}{m}/m_{\stackrel{~}{g}})`$ corrections, but neglecting terms $`𝒪[\alpha _{s,t}^{n+1}\mathrm{ln}^n(\stackrel{~}{m}/m_{\stackrel{~}{g}})]`$ with $`n0`$. For the radiative 2–body gluino decay into a gluon and a neutralino, a greater accuracy is more appropriate. The expressions of $`C_7^{\stackrel{~}{B}}`$ and $`C_5^{\stackrel{~}{H}}`$ given in eqs. (10) and (21), together with leading–order anomalous dimensions and one–loop matrix elements \[see eq. (62) below\], allow us to determine the 2–body decay amplitude neglecting terms $`𝒪[\alpha _{s,t}^{n+1}\mathrm{ln}^n(\stackrel{~}{m}/m_{\stackrel{~}{g}})]`$ with $`n1`$. This means that we have resummed all large logarithms at the leading order in all cases, but our formulae for 2–body gluino decays contain also the complete $`𝒪(\alpha _{s,t})`$ terms, relevant when the logarithm is not large.
(iii) If $`\stackrel{~}{m}`$ is close to the GUT scale, in presence of gauge–coupling unification there is no solid justification for the approximation of computing $`\alpha _s`$ contributions to the anomalous dimensions, neglecting electroweak corrections. However, because of the large SU(3) coefficients, we consider our approximation to be fairly adequate, even for $`\stackrel{~}{m}`$ as large as $`10^{13}`$ GeV, which is the maximum value of $`\stackrel{~}{m}`$ consistent with the negative searches for anomalous heavy isotopes.
(iv) In eq. (20) we have included the contribution from the bottom Yukawa coupling $`h_b`$, since these coefficients are enhanced when $`\mathrm{tan}\beta `$ is large. There are no $`\mathrm{tan}\beta `$ enhancements in the evolution below $`\stackrel{~}{m}`$, and therefore our results are reliable for any value of $`\mathrm{tan}\beta `$.
(v) Split Supersymmetry is free from flavour problems, therefore our assumption that squark mass matrices are diagonal is unnecessary. On the other hand, a certain degree of mass degeneracy among squarks is required by gauge-coupling unification. In the results presented in sect. 4 we take for simplicity all squark masses to be equal.
## 3 Operator Renormalization
The renormalization–group flow for the Wilson coefficients is determined by the equations
$`\mu {\displaystyle \frac{d\stackrel{}{C}}{d\mu }}`$ $`=`$ $`\widehat{\gamma }^T(\alpha _s,\alpha _t)\stackrel{}{C}`$ (22)
$`\mu {\displaystyle \frac{d\alpha _s}{d\mu }}`$ $`=`$ $`\beta _s{\displaystyle \frac{\alpha _s^2}{2\pi }}`$ (23)
$`\mu {\displaystyle \frac{d\alpha _t}{d\mu }}`$ $`=`$ $`\beta _t{\displaystyle \frac{\alpha _t^2}{2\pi }}\beta _{st}{\displaystyle \frac{\alpha _s\alpha _t}{2\pi }},`$ (24)
where $`\mu `$ is the renormalization scale and, in Split Supersymmetry, we have $`\beta _s=5`$, $`\beta _t=9/2`$ and $`\beta _{st}=8`$. The anomalous–dimension matrix $`\widehat{\gamma }`$ can be expressed as
$$\widehat{\gamma }_{ij}=2b_{ij}\delta _{ij}\underset{f}{}a_f,$$
(25)
where $`b_{ij}`$ are extracted from the poles of the one–loop renormalization of the operators $`Q_i`$ ($`Q_ib_{ij}Q_j/ϵ+\mathrm{}`$). In eq. (25) the sum is over all fields entering the operator $`Q_i`$, and the field anomalous dimensions $`a_f`$ are given by
$$a_{q_L^k}=\frac{1}{4\pi }\left(\alpha _sC_F+\frac{\alpha _t}{2}\delta _{k3}\right),a_{u_R^k}=\frac{1}{4\pi }\left(\alpha _sC_F+\alpha _t\delta _{k3}\right),a_{d_R}=\frac{\alpha _sC_F}{4\pi },$$
(26)
$$a_{\stackrel{~}{g}}=\frac{\alpha _sN_c}{4\pi },a_h=\frac{\alpha _tN_c}{4\pi },a_g=\frac{\alpha _s}{4\pi }\left(N_c\frac{2}{3}N_f\right).$$
(27)
Here $`k`$ is a generation index, $`C_F=(N_c^21)/(2N_c)`$, $`N_c=3`$, $`N_f=6`$. Note that the gluon anomalous dimension $`a_g`$ (given here in the Feynman gauge) is different from the SM value because it includes the gluino contribution.
We find that the anomalous–dimension matrices of the $`B`$–ino operators in eqs. (2)–(8), of the $`W`$–ino operators in eqs. (11)–(12), and of the higgsino operators in eqs. (14)–(18) are respectively
$$\widehat{\gamma }^{(a)}=\frac{\alpha _s}{4\pi }\gamma _s^a+\frac{\alpha _t}{4\pi }\gamma _t^a+\frac{\sqrt{\alpha _s\alpha _t}}{4\pi }\gamma _{st}^a,a=\stackrel{~}{B},\stackrel{~}{W},\stackrel{~}{H}$$
(28)
$$\gamma _s^{\stackrel{~}{B}}=\frac{1}{3}\left(\begin{array}{ccccccc}89N_c& 8& 8& 8& 8& 8& 0\\ 4& 49N_c& 4& 4& 4& 4& 0\\ 4& 4& 49N_c& 4& 4& 4& 0\\ 4& 4& 4& 49N_c& 4& 4& 0\\ 2& 2& 2& 2& 29N_c& 2& 0\\ 2& 2& 2& 2& 2& 29N_c& 0\\ 0& 0& 0& 0& 0& 0& 2N_f18N_c\end{array}\right),$$
(29)
$$\gamma _t^{\stackrel{~}{B}}=\left(\begin{array}{ccccccc}0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 2& 0& 0\\ 0& 0& 0& 1& 2& 0& 0\\ 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0\end{array}\right),\gamma _{st}^{\stackrel{~}{B}}=0,$$
(30)
$$\gamma _s^{\stackrel{~}{W}}=\left(\begin{array}{cc}3N_c& 0\\ 0& 3N_c\end{array}\right),\gamma _t^{\stackrel{~}{W}}=\left(\begin{array}{cc}0& 0\\ 0& 1\end{array}\right),\gamma _{st}^{\stackrel{~}{W}}=0,$$
(31)
$$\gamma _s^{\stackrel{~}{H}}=\left(\begin{array}{ccccc}\frac{3}{N_c}& 0& 0& 0& 0\\ 0& 4N_c\frac{1}{N_c}& 0& 0& 0\\ 0& 0& \frac{3}{N_c}& 0& 0\\ 0& 0& 0& 4N_c\frac{1}{N_c}& 0\\ 0& 0& 0& 0& \frac{2}{3}N_f6N_c\end{array}\right),$$
(32)
$$\gamma _t^{\stackrel{~}{H}}=\frac{1}{2}\left(\begin{array}{ccccc}3& 0& 0& 0& 0\\ 0& 3& 0& 0& 0\\ 0& 0& 1& 0& 0\\ 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 2N_c\end{array}\right),\gamma _{st}^{\stackrel{~}{H}}=\left(\begin{array}{ccccc}0& 0& 0& 0& 0\\ 0& 0& 0& 0& 4\\ 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0\\ 0& 2& 0& 0& 0\end{array}\right).$$
(33)
For coefficients with only multiplicative renormalization (which is the case for $`C_7^{\stackrel{~}{B}}`$, $`C_{1,2}^{\stackrel{~}{W}}`$, $`C_{1,3,4}^{\stackrel{~}{H}}`$), eq. (22) can be easily integrated, with the result
$$C_i(\mu )=C_i(\stackrel{~}{m})\eta _s^{\left(\frac{\gamma _s}{2\beta _s}\frac{\beta _{st}\gamma _t}{2\beta _s\beta _t}\right)}\eta _t^{\frac{\gamma _t}{2\beta _t}}\mathrm{for}C_i=C_7^{\stackrel{~}{B}},C_{1,2}^{\stackrel{~}{W}},C_{1,3,4}^{\stackrel{~}{H}}.$$
(34)
We have defined
$`\eta _s{\displaystyle \frac{\alpha _s(\stackrel{~}{m})}{\alpha _s(\mu )}}`$ $`=`$ $`1+{\displaystyle \frac{\alpha _s(\stackrel{~}{m})}{2\pi }}\beta _s\mathrm{ln}{\displaystyle \frac{\mu }{\stackrel{~}{m}}},`$ (35)
$`\eta _t{\displaystyle \frac{\alpha _t(\stackrel{~}{m})}{\alpha _t(\mu )}}`$ $`=`$ $`\eta _s^{\frac{\beta _{st}}{\beta _s}}+{\displaystyle \frac{\alpha _t(\stackrel{~}{m})\beta _t}{\alpha _s(\stackrel{~}{m})\left(\beta _{st}\beta _s\right)}}\left(\eta _s^{\frac{\beta _{st}}{\beta _s}}\eta _s\right).`$ (36)
The evolution of the Wilson coefficients for the other $`B`$–ino operators involves operator mixing and the solution of eq. (22) is given by
$`C_i^{\stackrel{~}{B}}(\mu )`$ $`=`$ $`\eta _s^{\frac{9}{10}}\left[C_i^{\stackrel{~}{B}}(\stackrel{~}{m})+y\overline{C}(\stackrel{~}{m})\right]i=1,2,3,6`$ (37)
$`C_4^{\stackrel{~}{B}}(\mu )`$ $`=`$ $`\eta _s^{\frac{9}{10}}\left[(1+z)C_4^{\stackrel{~}{B}}(\stackrel{~}{m})zC_5^{\stackrel{~}{B}}(\stackrel{~}{m})+y\overline{C}(\stackrel{~}{m})\right],`$ (38)
$`C_5^{\stackrel{~}{B}}(\mu )`$ $`=`$ $`\eta _s^{\frac{9}{10}}\left[(1+2z)C_5^{\stackrel{~}{B}}(\stackrel{~}{m})2zC_4^{\stackrel{~}{B}}(\stackrel{~}{m})+y\overline{C}(\stackrel{~}{m})\right],`$ (39)
where $`\overline{C}=C_1^{\stackrel{~}{B}}/3+(C_2^{\stackrel{~}{B}}+C_3^{\stackrel{~}{B}}+C_4^{\stackrel{~}{B}})/6+(C_5^{\stackrel{~}{B}}+C_6^{\stackrel{~}{B}})/12`$, $`y=\eta _s^{\mathrm{\hspace{0.17em}4}/5}1`$, and $`z=(\eta _s^{\mathrm{\hspace{0.17em}8}/15}\eta _t^{1/3}1)/3`$. Because of the non–vanishing contribution from $`\gamma _{st}^{\stackrel{~}{H}}`$, the equations for $`C_2^{\stackrel{~}{H}}`$ and $`C_5^{\stackrel{~}{H}}`$ cannot be solved analytically. The numerical results for the renormalization coefficients $`\mathrm{\Delta }_{ij}`$, defined by
$$\left(\begin{array}{c}C_2^{\stackrel{~}{H}}(\mu )\\ C_5^{\stackrel{~}{H}}(\mu )\end{array}\right)=\left(\begin{array}{cc}\mathrm{\Delta }_{22}& \mathrm{\Delta }_{25}\\ \mathrm{\Delta }_{52}& \mathrm{\Delta }_{55}\end{array}\right)\left(\begin{array}{c}C_2^{\stackrel{~}{H}}(\stackrel{~}{m})\\ C_5^{\stackrel{~}{H}}(\stackrel{~}{m})\end{array}\right),$$
(40)
are shown in fig. 1 for a representative choice of $`\alpha _s(\stackrel{~}{m})`$ and $`\alpha _t(\stackrel{~}{m})`$. Despite the fact that the high–energy boundary condition on $`C_5^{\stackrel{~}{H}}`$, eq. (21), is suppressed by a loop factor, a sizeable value of $`C_5^{\stackrel{~}{H}}`$ can be generated through the mixing with $`C_2^{\stackrel{~}{H}}`$.
A computation of the $`𝒪(\alpha _s)`$ part of the anomalous dimensions, restricted to the four–fermion operators, has been given in the appendix of ref. . From the comparison with eq. (29) it appears that the authors of ref. have omitted the mixing among the $`B`$–ino operators induced by the penguin diagrams. Also, we disagree with ref. on the anomalous dimensions of the higgsino operators.
Once we have evolved the Wilson coefficients down to the renormalization scale at which we compute the gluino decay width, it is convenient to express the operators in terms of chargino and neutralino mass eigenstates. With the usual definitions for the chargino and neutralino mixing matrices $`U`$, $`V`$ and $`N`$, which we assume to be real, the $`B`$–ino, $`W`$–ino and higgsino spinors can be expressed as
$$\overline{\stackrel{~}{W}^+}=\overline{\chi _i^+}\left(U_{i1}P_L+V_{i1}P_R\right),\overline{\stackrel{~}{H}^+}=\overline{\chi _i^+}\left(U_{i2}P_L+V_{i2}P_R\right),$$
(41)
$$\overline{\stackrel{~}{B}}=\overline{\chi _i^0}N_{i1},\overline{\stackrel{~}{W}^3}=\overline{\chi _i^0}N_{i2},\overline{\stackrel{~}{H}^0}=\overline{\chi _i^0}\left(N_{i4}P_RN_{i3}P_L\right),$$
(42)
where $`P_L`$ and $`P_R`$ are the chiral projectors. In the basis of mass eigenstates, the effective Lagrangian becomes
$$=\frac{1}{\stackrel{~}{m}^2}\underset{j}{}C_j^{\chi _i^0}Q_j^{\chi _i^0}+\frac{1}{\stackrel{~}{m}^2}(\underset{j}{}C_j^{\chi _i^+}Q_j^{\chi _i^+}+\mathrm{h}.\mathrm{c}.).$$
(43)
The operators involving neutralinos and quarks and their corresponding Wilson coefficients are
$`Q_{1q_L,q_R}^{\chi _i^0}`$ $`=`$ $`\overline{\chi _i^0}\gamma ^\mu \gamma _5\stackrel{~}{g}^a{\displaystyle \underset{k=1}{\overset{2}{}}}\overline{q}_{L,R}^{(k)}\gamma _\mu T^aq_{L,R}^{(k)}(q=u,d),`$ (44)
$`Q_{2q_L,q_R}^{\chi _i^0}`$ $`=`$ $`\overline{\chi _i^0}\gamma ^\mu \gamma _5\stackrel{~}{g}^a\overline{q}_{L,R}\gamma _\mu T^aq_{L,R}(q=t,b),`$ (45)
$`Q_{3q_L,q_R}^{\chi _i^0}`$ $`=`$ $`\overline{\chi _i^0}_{R,L}\stackrel{~}{g}_{L,R}^a\overline{q}_{R,L}T^aq_{L,R}(q=t,b),`$ (46)
$`Q_{4q_L,q_R}^{\chi _i^0}`$ $`=`$ $`\overline{\chi _i^0}_{R,L}\sigma ^{\mu \nu }\gamma _5\stackrel{~}{g}_{L,R}^a\overline{q}_{R,L}\sigma _{\mu \nu }T^aq_{L,R}(q=t,b),`$ (47)
$`C_{1u_L}^{\chi _i^0}=C_1^{\stackrel{~}{B}}N_{i1}+C_1^{\stackrel{~}{W}}N_{i2},C_{1u_R}^{\chi _i^0}=C_2^{\stackrel{~}{B}}N_{i1},`$ (48)
$`C_{1d_L}^{\chi _i^0}=C_1^{\stackrel{~}{B}}N_{i1}C_1^{\stackrel{~}{W}}N_{i2},C_{1d_R}^{\chi _i^0}=C_3^{\stackrel{~}{B}}N_{i1},`$ (49)
$`C_{2t_L}^{\chi _i^0}=C_4^{\stackrel{~}{B}}N_{i1}+C_2^{\stackrel{~}{W}}N_{i2},C_{3t_L}^{\chi _i^0}=C_1^{\stackrel{~}{H}}N_{i4},C_{4t_L}^{\chi _i^0}=C_2^{\stackrel{~}{H}}N_{i4},`$ (50)
$`C_{2t_R}^{\chi _i^0}=C_5^{\stackrel{~}{B}}N_{i1},C_{3t_R}^{\chi _i^0}=C_1^{\stackrel{~}{H}}N_{i4},C_{4t_R}^{\chi _i^0}=C_2^{\stackrel{~}{H}}N_{i4},`$ (51)
$`C_{2b_L}^{\chi _i^0}=C_4^{\stackrel{~}{B}}N_{i1}C_2^{\stackrel{~}{W}}N_{i2}C_{3b_L}^{\chi _i^0}=C_3^{\stackrel{~}{H}}N_{i3},C_{4b_L}^{\chi _i^0}=C_4^{\stackrel{~}{H}}N_{i3},`$ (52)
$`C_{2b_R}^{\chi _i^0}=C_6^{\stackrel{~}{B}}N_{i1},C_{3b_R}^{\chi _i^0}=C_3^{\stackrel{~}{H}}N_{i3},C_{4b_R}^{\chi _i^0}=C_4^{\stackrel{~}{H}}N_{i3}.`$ (53)
The operators involving charginos and quarks and their Wilson coefficients are
$`Q_{1L,R}^{\chi _i^+}`$ $`=`$ $`\overline{\chi _i^+}_{L,R}\gamma ^\mu \stackrel{~}{g}_{L,R}^a{\displaystyle \underset{k=1}{\overset{2}{}}}\overline{d}_L^{(k)}\gamma _\mu T^au_L^{(k)}`$ (54)
$`Q_{2L,R}^{\chi _i^+}`$ $`=`$ $`\overline{\chi _i^+}_{L,R}\gamma ^\mu \stackrel{~}{g}_{L,R}^a\overline{b}_L\gamma _\mu T^at_L`$ (55)
$`Q_{3L,R}^{\chi _i^+}`$ $`=`$ $`\overline{\chi _i^+}_{R,L}\stackrel{~}{g}_{L,R}^a\overline{b}_{R,L}T^at_{L,R}`$ (56)
$`Q_{4L,R}^{\chi _i^+}`$ $`=`$ $`\overline{\chi _i^+}_{R,L}\sigma ^{\mu \nu }\stackrel{~}{g}_{L,R}^a\overline{b}_{R,L}\sigma _{\mu \nu }T^at_{L,R}`$ (57)
$`C_{1L}^{\chi _i^+}=\sqrt{2}C_1^{\stackrel{~}{W}}V_{i1},C_{1R}^{\chi _i^+}=\sqrt{2}C_1^{\stackrel{~}{W}}U_{i1},`$ (58)
$`C_{2L}^{\chi _i^+}=\sqrt{2}C_2^{\stackrel{~}{W}}V_{i1},C_{3L}^{\chi _i^+}=C_3^{\stackrel{~}{H}}U_{i2},C_{4L}^{\chi _i^+}=C_4^{\stackrel{~}{H}}U_{i2},`$ (59)
$`C_{2R}^{\chi _i^+}=\sqrt{2}C_2^{\stackrel{~}{W}}U_{i1},C_{3R}^{\chi _i^+}=C_1^{\stackrel{~}{H}}V_{i2},C_{4R}^{\chi _i^+}=C_2^{\stackrel{~}{H}}V_{i2}.`$ (60)
All Wilson coefficients in eqs.(48)–(53) and (58)–(60) are evaluated at the scale $`\mu `$ at which the gluino decay width is computed (we take $`\mu =m_{\stackrel{~}{g}}`$ in our numerical analysis).
The magnetic operator involving a neutralino and a gluon is
$$Q_g^{\chi _i^0}=\overline{\chi _i^0}\sigma ^{\mu \nu }\gamma _5\stackrel{~}{g}^aG_{\mu \nu }^a.$$
(61)
In order to reach the desired accuracy in the $`\stackrel{~}{g}g\stackrel{~}{\chi }^0`$ process, we need to include the matrix element contribution coming from the diagram in which the two top quarks in the operator $`Q_2^{\stackrel{~}{H}}`$ close in a loop emitting a gluon. This results into an “effective” Wilson coefficient
$$C_{g}^{\chi _i^0}{}_{\mathrm{eff}}{}^{}(\mu )=C_7^{\stackrel{~}{B}}(\mu )N_{i1}+C_5^{\stackrel{~}{H}}(\mu )N_{i4}v+\frac{g_sh_t}{8\pi ^2}C_2^{\stackrel{~}{H}}(\mu )N_{i4}v\mathrm{ln}\frac{m_t^2}{\mu ^2},$$
(62)
where $`v`$ is the Higgs vacuum expectation value and we take $`\mu =m_{\stackrel{~}{g}}`$.
From the effective Lagrangian of eq. (43) we can compute the gluino decay widths and complete expressions can be found in the appendix. The same effective Lagrangian correctly describes also the interactions that lead to the decays $`\stackrel{~}{g}gg\stackrel{~}{\chi }^0`$ and $`\stackrel{~}{g}gh^0\stackrel{~}{\chi }^0`$. However, since these processes are subleading, we will not explicitly calculate their decay widths.
## 4 Results
We are now ready to discuss the results of our computation of the decay width and branching ratios of the gluino in Split Supersymmetry. The input parameters relevant to our analysis are the sfermion mass scale $`\stackrel{~}{m}`$, the physical gluino mass $`m_{\stackrel{~}{g}}`$ and $`\mathrm{tan}\beta `$, which in Split Supersymmetry is interpreted as the tangent of the angle that rotates the finely tuned Higgs doublets. To simplify the analysis we assume that the squark masses are degenerate, i.e. we set $`r_{\stackrel{~}{q}_L}=r_{\stackrel{~}{u}_R}=r_{\stackrel{~}{d}_R}=1`$ in the matching conditions of the Wilson coefficients. The gluino mass parameter in the Lagrangian, $`M_3`$, is extracted from $`m_{\stackrel{~}{g}}`$ including radiative corrections, and the other gaugino masses $`M_1`$ and $`M_2`$ are computed from $`M_3`$ assuming unification at the GUT scale. The higgsino mass parameter $`\mu `$ is determined as a function of $`M_2`$ by requiring that the relic abundance of neutralinos is equal to the dark–matter density preferred by WMAP data (see fig. 11 of ref. ). The sign of $`\mu `$ remains a free parameter, but since it does not affect our results for the gluino decays in a significant way we will assume $`\mu >0`$ throughout our analysis. The effective couplings of gauginos and higgsinos at the weak scale, needed to compute the chargino and neutralino mass matrices, are determined from their high–energy (supersymmetric) boundary values by means of the renormalization–group equations of Split Supersymmetry, given in ref. . Finally, the SM input parameters relevant to our analysis are: the physical masses for the top quark and gauge bosons, $`m_t=178`$ GeV, $`M_Z=91.187`$ GeV and $`M_W=80.41`$ GeV; the running bottom mass computed at the scale of the top mass, $`m_b(m_t)=2.75`$ GeV; the Fermi constant, $`G_F=1.166\times 10^5`$ GeV<sup>-2</sup>; the running strong coupling computed at the scale of the top mass, $`\alpha _s(m_t)=0.106`$.
To start our discussion, we show in fig. 2 the gluino lifetime $`\tau _{\stackrel{~}{g}}`$ (in seconds) as a function of the sfermion mass scale $`\stackrel{~}{m}`$, for $`\mathrm{tan}\beta =2`$ and four different values of the physical gluino mass ($`m_{\stackrel{~}{g}}`$ = 0.5, 1, 2 and 5 TeV, respectively). It can be seen that $`\tau _{\stackrel{~}{g}}`$ is about 4 seconds for $`m_{\stackrel{~}{g}}=1`$ TeV and $`\stackrel{~}{m}=10^9`$ GeV. A value of $`\tau _{\stackrel{~}{g}}`$ equal to the age of the universe (14 Gyr) corresponds to $`\stackrel{~}{m}=(1.1,\mathrm{\hspace{0.17em}2.1},\mathrm{\hspace{0.17em}4.5},\mathrm{\hspace{0.17em}13})\times 10^{13}`$ GeV for $`m_{\stackrel{~}{g}}`$ = 0.5, 1, 2 and 5 TeV, respectively.
In fig. 3 we show the branching ratios for the three decay processes $`\stackrel{~}{g}\chi ^0g`$, $`\stackrel{~}{g}\chi ^0q\overline{q}`$ and $`\stackrel{~}{g}\chi ^\pm q\overline{q}^{}`$ (summed over all neutralino or chargino states) as a function of $`\stackrel{~}{m}`$, for $`\mathrm{tan}\beta =20`$ and three different values of $`m_{\stackrel{~}{g}}`$ : 500 GeV (upper plots), 1 TeV (middle plots) and 2 TeV (lower plots). The value of $`\mathrm{tan}\beta `$ has little impact on these results. The plots on the left of fig. 3 represent the results of our full calculation, including the resummation of the leading logarithmic corrections controlled by $`\alpha _s`$ and $`\alpha _t`$. The plots on the right represent instead the lowest–order results that do not include the resummation. We obtain the latter results by replacing the Wilson coefficients of the four–fermion operators in the low–energy effective Lagrangian with their tree–level expressions in terms of gauge and Yukawa couplings \[eqs. (9)–(10), (13) and (19)–(20)\], and the Wilson coefficient of the magnetic operator with its one–loop expression. The plots in fig. 3 show that the branching ratio of the radiative decay $`\stackrel{~}{g}\chi ^0g`$ decreases for increasing $`m_{\stackrel{~}{g}}`$ and increases for increasing $`\stackrel{~}{m}`$. In fact, as stressed in ref. , the ratio between the two–body and three–body decay rates computed at lowest order scales like $`m_t^2/m_{\stackrel{~}{g}}^2[1\mathrm{ln}(\stackrel{~}{m}^2/m_t^2)]^2`$, where the logarithmic term comes from the top–stop loop that generates the magnetic gluino–gluon–higgsino interaction. For large values of $`\stackrel{~}{m}`$, the resummation of the logarithms becomes necessary. Comparing the plots on the left and right sides of fig. 3, we see that the resummation of the leading logarithmic corrections tends to enhance the three–body decays and suppress the radiative decay. The effect of the corrections on the branching ratios is particularly visible when, like in the middle and lower plots, neither the two–body nor the three–body channels are obviously dominant in the range $`10^8`$ GeV $`<\stackrel{~}{m}<10^{13}`$ GeV, relevant to Split Supersymmetry.
To further illustrate the effect of the resummation of the leading logarithmic corrections, we plot in fig. 4 the ratio $`\mathrm{\Gamma }/\mathrm{\Gamma }_0`$ of the partial decay widths with and without resummation, for the processes $`\stackrel{~}{g}\chi ^0g`$, $`\stackrel{~}{g}\chi ^0q\overline{q}`$ and $`\stackrel{~}{g}\chi ^\pm q\overline{q}^{}`$. We fix $`m_{\stackrel{~}{g}}=1`$ TeV, $`\mathrm{tan}\beta =20`$ and $`\mu >0`$, but we have checked that the qualitative behaviour of the corrections is independent of the precise choice of the parameters. It can be seen from fig. 4 that for large enough values of $`\stackrel{~}{m}`$ the radiative corrections can be of the order of 50–100%, and that they enhance the widths for the three–body decays and suppress the width for the radiative decay.
To conclude this section, we discuss the scaling behaviour of the gluino lifetime and total decay width. The lifetime $`\tau _{\stackrel{~}{g}}=\mathrm{}/\mathrm{\Gamma }_{\mathrm{tot}}`$ can be written as
$$\tau _{\stackrel{~}{g}}=\frac{4\mathrm{sec}}{N}\times \left(\frac{\stackrel{~}{m}}{10^9\mathrm{GeV}}\right)^4\times \left(\frac{1\mathrm{TeV}}{m_{\stackrel{~}{g}}}\right)^5,$$
(63)
where the normalization $`N`$ is of order unity and depends on $`\stackrel{~}{m}`$ and $`m_{\stackrel{~}{g}}`$ (and only very mildly on $`\mathrm{tan}\beta `$). In fig. 5 we show $`N`$ as a function of $`\stackrel{~}{m}`$ for $`\mathrm{tan}\beta =20`$ and three different values of the physical gluino mass ($`m_{\stackrel{~}{g}}=0.5,\mathrm{\hspace{0.17em}1}`$ and 2 TeV, respectively). The non–vanishing slope of $`N`$ represents the deviation of the total gluino decay width from the naive scaling behaviour $`\mathrm{\Gamma }_{\mathrm{tot}}m_{\stackrel{~}{g}}^5/\stackrel{~}{m}^4`$. The solid lines in the plot represent the results of our full calculation, whereas the dashed lines represent the lowest–order results that do not include the resummation. For low values of $`m_{\stackrel{~}{g}}`$ the contribution of the radiative decay dominates (see fig. 3), thus the total decay width departs visibly from the naive scaling and is significantly suppressed by the resummation of the radiative corrections. On the other hand, for large values of $`m_{\stackrel{~}{g}}`$ the three–body decays dominate, and the effect of the resummation is to enhance the total decay width. Finally, for the intermediate value $`m_{\stackrel{~}{g}}=1`$ TeV there is a compensation between the corrections to the radiative decay width and those to the three–body decay widths, and the net effect on the total decay width of the resummation of the leading logarithmic corrections is rather small.
## 5 Gluino Decay into Gravitinos
Split Supersymmetry opens up the possibility of direct tree-level mediation of the original supersymmetry breaking to the SM superfields, without the need of a hidden sector . In usual low-energy supersymmetry, this possibility is impracticable: for $`F`$–term breaking some scalars remain lighter than the SM matter fermions, and for $`D`$–term breaking gaugino masses cannot be generated at the same order of scalar masses. In Split Supersymmetry a large hierarchy between scalar and gaugino masses is acceptable, and indeed models have been proposed with direct mediation of $`D`$–term supersymmetry breaking.
Therefore, in Split Supersymmetry the original scale of supersymmetry breaking $`\sqrt{F}`$, which is related to the gravitino mass by
$$m_{3/2}=\sqrt{\frac{8\pi }{3}}\frac{F}{M_{\mathrm{Pl}}},$$
(64)
could be as low as the squark mass scale $`\stackrel{~}{m}`$. This means that the interactions between the gluino and (the spin–1/2 component of) the gravitino, which are suppressed by $`1/F`$, could be as strong as those considered in the previous sections, which are suppressed by $`1/\stackrel{~}{m}^2`$.
For $`m_{3/2}m_{\stackrel{~}{g}}`$, the gravitino interactions can be obtained, through the supersymmetric analogue of the equivalence theorem , from the goldstino derivative coupling to the supercurrent. This approximation is valid as long as $`\sqrt{F}6\times (m_{\stackrel{~}{g}}/1\mathrm{TeV})^{1/2}\times 10^{10}`$ GeV. Using the equations of motion, we can write the effective goldstino ($`\stackrel{~}{G}`$) interactions for on–shell particles as
$$=\frac{1}{F}\left(m_{\stackrel{~}{q}_L}^2\stackrel{~}{q}_L\overline{q}_Lm_{\stackrel{~}{q}_R}^2\stackrel{~}{q}_R\overline{q}_R+\frac{m_{\stackrel{~}{g}}}{4\sqrt{2}}\overline{\stackrel{~}{g}^a}\sigma ^{\mu \nu }\gamma _5G_{\mu \nu }^a\right)\stackrel{~}{G}+\mathrm{h}.\mathrm{c}.$$
(65)
Below $`\stackrel{~}{m}`$, the effective Lagrangian describing the interactions between the gluino and the goldstino becomes
$$=\frac{1}{F}\underset{i=1}{\overset{5}{}}C_i^{\stackrel{~}{G}}Q_i^{\stackrel{~}{G}},$$
(66)
$`Q_1^{\stackrel{~}{G}}`$ $`=`$ $`\overline{\stackrel{~}{G}}\gamma ^\mu \gamma _5\stackrel{~}{g}^a{\displaystyle \underset{\genfrac{}{}{0pt}{}{k=1,2}{q=u,d}}{}}\overline{q}^{(k)}\gamma _\mu T^aq^{(k)}`$ (67)
$`Q_2^{\stackrel{~}{G}}`$ $`=`$ $`\overline{\stackrel{~}{G}}\gamma ^\mu \gamma _5\stackrel{~}{g}^a\overline{q}_L^{(3)}\gamma _\mu T^aq_L^{(3)}`$ (68)
$`Q_3^{\stackrel{~}{G}}`$ $`=`$ $`\overline{\stackrel{~}{G}}\gamma ^\mu \gamma _5\stackrel{~}{g}^a\overline{t}_R\gamma _\mu T^at_R`$ (69)
$`Q_4^{\stackrel{~}{G}}`$ $`=`$ $`\overline{\stackrel{~}{G}}\gamma ^\mu \gamma _5\stackrel{~}{g}^a\overline{b}_R\gamma _\mu T^ab_R`$ (70)
$`Q_5^{\stackrel{~}{G}}`$ $`=`$ $`\overline{\stackrel{~}{G}}\sigma ^{\mu \nu }\gamma _5\stackrel{~}{g}^aG_{\mu \nu }^a.`$ (71)
The Wilson coefficients at the matching scale $`\stackrel{~}{m}`$ are
$$C_1^{\stackrel{~}{G}}=C_2^{\stackrel{~}{G}}=C_3^{\stackrel{~}{G}}=C_4^{\stackrel{~}{G}}=\frac{g_s}{\sqrt{2}},C_5^{\stackrel{~}{G}}=\frac{m_{\stackrel{~}{g}}}{2\sqrt{2}}.$$
(72)
Note that the coefficients of the interactions in eq. (66) have no dependence on $`\stackrel{~}{m}`$, because the squark mass square in the propagators of the particles we integrate out is exactly cancelled by the squark mass square in the goldstino coupling in eq. (65).
The operator renormalization proceeds analogously to the discussion in sect. 3. The anomalous dimension matrix of the operators in eqs. (67)–(71) is given by eq. (28) with
$$\gamma _s^{\stackrel{~}{G}}=\frac{1}{3}\left(\begin{array}{ccccc}169N_c& 16& 16& 16& 0\\ 4& 49N_c& 4& 4& 0\\ 2& 2& 29N_c& 2& 0\\ 2& 2& 2& 29N_c& 0\\ 0& 0& 0& 0& 2N_f18N_c\end{array}\right),$$
(73)
$$\gamma _t^{\stackrel{~}{G}}=\left(\begin{array}{ccccc}0& 0& 0& 0& 0\\ 0& 1& 2& 0& 0\\ 0& 1& 2& 0& 0\\ 0& 0& 0& 0& 0\end{array}\right),\gamma _{st}^{\stackrel{~}{G}}=0.$$
(74)
The evolution of the Wilson coefficients for the goldstino operators has the simple analytic form
$`C_i^{\stackrel{~}{G}}(\mu )`$ $`=`$ $`\eta _s^{\frac{9}{10}}\left[C_i^{\stackrel{~}{G}}(\stackrel{~}{m})+y\overline{C}^{\stackrel{~}{G}}(\stackrel{~}{m})\right]i=1,4`$ (75)
$`C_2^{\stackrel{~}{G}}(\mu )`$ $`=`$ $`\eta _s^{\frac{9}{10}}\left[(1+z)C_2^{\stackrel{~}{G}}(\stackrel{~}{m})zC_3^{\stackrel{~}{G}}(\stackrel{~}{m})+y\overline{C}^{\stackrel{~}{G}}(\stackrel{~}{m})\right],`$ (76)
$`C_3^{\stackrel{~}{G}}(\mu )`$ $`=`$ $`\eta _s^{\frac{9}{10}}\left[(1+2z)C_3^{\stackrel{~}{G}}(\stackrel{~}{m})2zC_2^{\stackrel{~}{G}}(\stackrel{~}{m})+y\overline{C}^{\stackrel{~}{G}}(\stackrel{~}{m})\right],`$ (77)
$`C_5^{\stackrel{~}{G}}(\mu )`$ $`=`$ $`\eta _s^{\frac{7}{5}}C_5^{\stackrel{~}{G}}(\stackrel{~}{m}),`$ (78)
where $`\overline{C}^{\stackrel{~}{G}}=2C_1^{\stackrel{~}{G}}/3+C_2^{\stackrel{~}{G}}/6+(C_3^{\stackrel{~}{G}}+C_4^{\stackrel{~}{G}})/12`$, $`y=\eta _s^{\mathrm{\hspace{0.17em}4}/5}1`$, and $`z=(\eta _s^{\mathrm{\hspace{0.17em}8}/15}\eta _t^{1/3}1)/3`$. The quantities $`\eta _s`$ and $`\eta _t`$ have been defined in eqs. (35) and (36), respectively.
The formulae for the gluino decay widths into goldstino and quarks and into goldstino and gluon can be found in the appendix. In fig. 6 we show the branching ratio for the process $`\stackrel{~}{g}\stackrel{~}{G}g`$ as a function of the ratio $`\sqrt{F}/\stackrel{~}{m}`$, for $`\stackrel{~}{m}=10^9`$ GeV, $`\mathrm{tan}\beta =2,\mu >0`$ and different values of the gluino mass. The branching ratio for the decay into goldstino and quarks, suppressed by phase space, is always at or below the 1% level. It can be seen from fig. 6 that the gluino decay into goldstino and gluon is largely dominant when $`\sqrt{F}`$ is as low as $`\stackrel{~}{m}`$. In fact, the decays into charginos or neutralinos and quarks (relevant for large values of $`m_{\stackrel{~}{g}}`$) are suppressed by phase space, while the radiative decay into gluon and neutralinos (relevant for smaller values of $`m_{\stackrel{~}{g}}`$) is suppressed by $`m_t^2/m_{\stackrel{~}{g}}^2`$ and a loop factor. With respect to the scaling behaviour outlined in eq. (63), the additional contribution to the total gluino decay width coming from the decay into goldstino and gluon can significantly suppress the gluino lifetime. In fact, for $`\sqrt{F}=\stackrel{~}{m}`$ the normalization $`N`$ in eq. (63) takes on values of order 40–50 for $`\stackrel{~}{m}>10^8`$ GeV.
On the other hand, the widths for the gluino decays into goldstino are suppressed by a factor $`\stackrel{~}{m}^4/F^2`$ with respect to those for decays into charginos or neutralinos. Fig. 6 shows that as soon as we depart from the condition $`\sqrt{F}=\stackrel{~}{m}`$ the branching ratio for $`\stackrel{~}{g}\stackrel{~}{G}g`$ falls off very quickly, and already for $`\sqrt{F}/\stackrel{~}{m}`$ as large as 10 the gluino decays into goldstino are below the 1% level.
## 6 Conclusions
If Split Supersymmetry is the correct theory to describe physics beyond the Standard Model, one of its most spectacular manifestations might be the discovery of a very long–lived gluino at the LHC. In this paper we provided a precise determination of the gluino lifetime and branching ratios. Applying to Split Supersymmetry the effective Lagrangian and renormalization group techniques, we discussed the proper treatment of the radiative corrections that are enhanced by the large logarithm of the ratio between the sfermion mass scale and the gluino mass. We computed the anomalous dimensions of the operators relevant to the gluino decay, that allow us to resum to all orders in the perturbative expansion the leading logarithmic corrections controlled by $`\alpha _s`$ and $`\alpha _t`$. We also provided explicit analytical formulae for the gluino decay widths in terms of the Wilson coefficients of the effective Lagrangian of Split Supersymmetry. For representative values of the input parameters, we discussed the numerical impact of the radiative corrections and found that they can modify substantially the gluino decay width and branching ratios. Finally, we considered models with direct mediation of supersymmetry breaking, and we found that the gluino decays into gravitinos might dominate over the other decay modes.
## Appendix
We present in this appendix the explicit formulae for the leading three–body and two–body gluino decay widths. All the results are expressed in terms of the Wilson coefficients of the effective Lagrangian of Split Supersymmetry, discussed in sects. 2, 3 and 5.
### Three–body decays into quarks and chargino or neutralino:
denoting the momenta of the decay products as $`(p_1,p_2,p_3)(p_{q_I},p_{\overline{q}_J},p_\chi )`$, and $`s_{ij}=(p_i+p_j)^2`$, the three–body decay amplitude is given by
$$\mathrm{\Gamma }_{\chi q_I\overline{q}_J}=\frac{1}{256\pi ^3m_{\stackrel{~}{g}}^3\stackrel{~}{m}^4}\overline{\left|\right|^2}𝑑s_{13}𝑑s_{23}.$$
(79)
The bar over $`\left|\right|^2`$ denotes the average over colour and spin of the gluino and the sum over colour and spin of the final state particles (the dependence on $`\stackrel{~}{m}`$ has been factored out). The limits of the integration in the ($`s_{13},s_{23}`$) plane are
$`s_{13}^{\mathrm{max}}`$ $`=`$ $`m_{q_I}^2+m_\chi ^2+{\displaystyle \frac{1}{2s_{23}}}[(m_{\stackrel{~}{g}}^2m_{q_I}^2s_{23})(s_{23}m_{\overline{q}_J}^2+m_\chi ^2)`$ (80)
$`+\lambda ^{1/2}(s_{23},m_{\stackrel{~}{g}}^2,m_{q_I}^2)\lambda ^{1/2}(s_{23},m_{\overline{q}_J}^2,m_\chi ^2)],`$
$`s_{13}^{\mathrm{min}}`$ $`=`$ $`m_{q_I}^2+m_\chi ^2+{\displaystyle \frac{1}{2s_{23}}}[(m_{\stackrel{~}{g}}^2m_{q_I}^2s_{23})(s_{23}m_{\overline{q}_J}^2+m_\chi ^2)`$ (81)
$`\lambda ^{1/2}(s_{23},m_{\stackrel{~}{g}}^2,m_{q_I}^2)\lambda ^{1/2}(s_{23},m_{\overline{q}_J}^2,m_\chi ^2)],`$
$`s_{23}^{\mathrm{max}}`$ $`=`$ $`(|m_{\stackrel{~}{g}}|m_{q_I})^2,`$ (82)
$`s_{23}^{\mathrm{min}}`$ $`=`$ $`(|m_\chi |+m_{\overline{q}_J})^2,`$ (83)
where $`\lambda (x,y,z)=x^2+y^2+z^22(xy+xz+yz)`$.
In the computation of the decays involving quarks of the first and second generation we can neglect the quark masses and we find
$`\mathrm{\Gamma }_{\chi _i^+d\overline{u}}=\mathrm{\Gamma }_{\chi _i^{}u\overline{d}}`$ $`=`$ $`{\displaystyle \frac{m_{\stackrel{~}{g}}^5}{1536\pi ^3\stackrel{~}{m}^4}}\left[\left(C_{1L}^{\chi _i^+}{}_{}{}^{2}+C_{1R}^{\chi _i^+}{}_{}{}^{2}\right)g(x_i)2C_{1L}^{\chi _i^+}C_{1R}^{\chi _i^+}f(x_i)\right],`$ (84)
$`\mathrm{\Gamma }_{\chi _i^0q\overline{q}}`$ $`=`$ $`{\displaystyle \frac{m_{\stackrel{~}{g}}^5}{768\pi ^3\stackrel{~}{m}^4}}(C_{1q_L}^{\chi _i^0}{}_{}{}^{2}+C_{1q_R}^{\chi _i^0}{}_{}{}^{2})[g(x_i)+f(x_i)](q=u,d),`$ (85)
where $`x_i=m_{\chi _i}/m_{\stackrel{~}{g}}`$, and we have included an overall factor 2 to take into account the sum over the two generations of light quarks. The functions $`f`$ and $`g`$ are defined as:
$`g(x)`$ $`=`$ $`18x^2+8x^6x^812x^4\mathrm{ln}x^2,`$ (86)
$`f(x)`$ $`=`$ $`2x+18x^318x^52x^7+12x^3(1+x^2)\mathrm{ln}x^2.`$ (87)
For generic quark masses the integration of the squared amplitude $`\overline{\left|\right|^2}`$ on the $`(s_{13},s_{23})`$ plane cannot be performed analytically, and in order to compute the total decay width we must resort to a numerical integration.
The squared amplitude for the processes $`\stackrel{~}{g}\chi _i^+b\overline{t}`$ and $`\stackrel{~}{g}\chi _i^{}t\overline{b}`$ is given by
$`\overline{\left|\right|^2}`$ $`=`$ $`C_{2L}^{\chi _i^+}{}_{}{}^{2}(m_{\stackrel{~}{g}}^2+m_t^2s_{13})(s_{13}m_{\chi _i^+}^2m_b^2)`$ (88)
$`+`$ $`C_{2R}^{\chi _i^+}{}_{}{}^{2}(m_{\stackrel{~}{g}}^2+m_b^2s_{23})(s_{23}m_{\chi _i^+}^2m_t^2)`$
$`+`$ $`{\displaystyle \frac{1}{4}}\left(C_{3L}^{\chi _i^+}{}_{}{}^{2}+C_{3R}^{\chi _i^+}{}_{}{}^{2}\right)(m_{\chi _i^+}^2+m_{\stackrel{~}{g}}^2s_{13}s_{23})(s_{13}+s_{23}m_t^2m_b^2)`$
$`+`$ $`4(C_{4L}^{\chi _i^+}{}_{}{}^{2}+C_{4R}^{\chi _i^+}{}_{}{}^{2})[(m_{\chi _i^+}^2+m_{\stackrel{~}{g}}^2s_{13}s_{23})(s_{13}+s_{23}m_t^2m_b^24m_{\chi _i^+}^2)`$
$`+4(s_{13}m_{\chi _i^+}^2)(s_{23}m_{\chi _i^+}^2)4m_t^2m_b^2]`$
$`+`$ $`2C_{2L}^{\chi _i^+}C_{2R}^{\chi _i^+}m_{\stackrel{~}{g}}m_{\chi _i^+}(s_{13}+s_{23}m_{\chi _i^+}^2m_{\stackrel{~}{g}}^2)`$
$`+`$ $`\left(C_{2R}^{\chi _i^+}C_{3R}^{\chi _i^+}+12C_{2R}^{\chi _i^+}C_{4R}^{\chi _i^+}\right)m_{\chi _i^+}m_t(s_{23}m_b^2m_{\stackrel{~}{g}}^2)`$
$`+`$ $`\left(C_{2R}^{\chi _i^+}C_{3L}^{\chi _i^+}+12C_{2R}^{\chi _i^+}C_{4L}^{\chi _i^+}\right)m_{\stackrel{~}{g}}m_b(s_{23}m_t^2m_{\chi _i^+}^2)`$
$``$ $`\left(C_{2L}^{\chi _i^+}C_{3R}^{\chi _i^+}12C_{2L}^{\chi _i^+}C_{4R}^{\chi _i^+}\right)m_{\stackrel{~}{g}}m_t(s_{13}m_b^2m_{\chi _i^+}^2)`$
$``$ $`\left(C_{2L}^{\chi _i^+}C_{3L}^{\chi _i^+}12C_{2L}^{\chi _i^+}C_{4L}^{\chi _i^+}\right)m_{\chi _i^+}m_b(s_{13}m_t^2m_{\stackrel{~}{g}}^2)`$
$`+`$ $`2(C_{3L}^{\chi _i^+}C_{4L}^{\chi _i^+}+C_{3R}^{\chi _i^+}C_{4R}^{\chi _i^+})[(m_{\stackrel{~}{g}}^2+m_{\chi _i^+}^2s_{13}s_{23})(s_{23}s_{13}+m_b^2m_t^2)`$
$`+2m_b^2(s_{23}m_t^2m_{\chi _i^+}^2)2m_t^2(s_{13}m_b^2m_{\chi _i^+}^2)]`$
$``$ $`2\left(C_{3L}^{\chi _i^+}C_{3R}^{\chi _i^+}+48C_{4L}^{\chi _i^+}C_{4R}^{\chi _i^+}\right)m_{\stackrel{~}{g}}m_{\chi _i^+}m_tm_b.`$
The squared amplitude for the processes $`\stackrel{~}{g}\chi _i^0t\overline{t}`$ and $`\stackrel{~}{g}\chi _i^0b\overline{b}`$ is given by
$`\overline{\left|\right|^2}`$ $`=`$ $`(C_{2q_L}^{\chi _i^0}{}_{}{}^{2}+C_{2q_R}^{\chi _i^0}{}_{}{}^{2})[(m_{\stackrel{~}{g}}^2+m_q^2s_{13})(s_{13}m_q^2m_{\chi _i^0}^2)`$ (89)
$`+(m_{\stackrel{~}{g}}^2+m_q^2s_{23})(s_{23}m_q^2m_{\chi _i^0}^2)+2m_{\stackrel{~}{g}}m_{\chi _i^0}(m_{\stackrel{~}{g}}^2+m_{\chi _i^0}^2s_{13}s_{23})]`$
$`+`$ $`{\displaystyle \frac{1}{4}}\left(C_{3q_L}^{\chi _i^0}{}_{}{}^{2}+C_{3q_R}^{\chi _i^0}{}_{}{}^{2}\right)(m_{\chi _i^0}^2+m_{\stackrel{~}{g}}^2s_{13}s_{23})(s_{13}+s_{23}2m_q^2)`$
$`+`$ $`4(C_{4q_L}^{\chi _i^0}{}_{}{}^{2}+C_{4q_R}^{\chi _i^0}{}_{}{}^{2})[(m_{\chi _i^0}^2+m_{\stackrel{~}{g}}^2+4m_q^2s_{13}s_{23})(s_{13}+s_{23}2m_q^24m_{\chi _i^0}^2)`$
$`+4(s_{13}m_q^2m_{\chi _i^0}^2)(s_{23}m_q^2m_{\chi _i^0}^2)+8m_q^2m_{\chi _i^0}^2)]`$
$`+`$ $`4C_{2q_L}^{\chi _i^0}C_{2q_R}^{\chi _i^0}m_q^2(s_{13}+s_{23}+4m_{\stackrel{~}{g}}m_{\chi _i^0}2m_q^2)`$
$`+`$ $`\left(C_{2q_R}^{\chi _i^0}C_{3q_R}^{\chi _i^0}C_{2q_L}^{\chi _i^0}C_{3q_L}^{\chi _i^0}\right)m_q\left[m_{\chi _i^0}(m_{\stackrel{~}{g}}^2+m_q^2s_{13})+m_{\stackrel{~}{g}}(m_{\chi _i^0}^2+m_q^2s_{23})\right]`$
$`+`$ $`\left(C_{2q_R}^{\chi _i^0}C_{3q_L}^{\chi _i^0}C_{2q_L}^{\chi _i^0}C_{3q_R}^{\chi _i^0}\right)m_q\left[m_{\chi _i^0}(m_{\stackrel{~}{g}}^2+m_q^2s_{23})+m_{\stackrel{~}{g}}(m_{\chi _i^0}^2+m_q^2s_{13})\right]`$
$`+`$ $`12\left(C_{2q_R}^{\chi _i^0}C_{4q_R}^{\chi _i^0}C_{2q_L}^{\chi _i^0}C_{4q_L}^{\chi _i^0}\right)m_q\left[m_{\stackrel{~}{g}}(m_{\chi _i^0}^2+m_q^2s_{23})m_{\chi _i^0}(m_{\stackrel{~}{g}}^2+m_q^2s_{13})\right]`$
$`+`$ $`12\left(C_{2q_R}^{\chi _i^0}C_{4q_L}^{\chi _i^0}C_{2q_L}^{\chi _i^0}C_{4q_R}^{\chi _i^0}\right)m_q\left[m_{\chi _i^0}(m_{\stackrel{~}{g}}^2+m_q^2s_{13})m_{\stackrel{~}{g}}(m_{\chi _i^0}^2+m_q^2s_{23})\right]`$
$`+`$ $`2\left(C_{3q_L}^{\chi _i^0}C_{4q_L}^{\chi _i^0}+C_{3q_R}^{\chi _i^0}C_{4q_R}^{\chi _i^0}\right)(m_{\stackrel{~}{g}}^2+m_{\chi _i^0}^2+2m_q^2s_{13}s_{23})(s_{23}s_{13})`$
$``$ $`2\left(C_{3q_L}^{\chi _i^0}C_{3q_R}^{\chi _i^0}+48C_{4q_L}^{\chi _i^0}C_{4q_R}^{\chi _i^0}\right)m_{\stackrel{~}{g}}m_{\chi _i^0}m_q^2(q=t,b).`$
We have checked that inserting in eqs. (88)–(89) the high–energy (i.e. non resummed) expressions for the Wilson coefficients given in sects. 2 and 3 we reproduce the tree–level results of ref. .
### Two–body decays into neutralino and gluon:
the width for the radiative decay of the gluino, $`\stackrel{~}{g}g\chi _i^0`$, is
$$\mathrm{\Gamma }_{\chi _i^0g}=\frac{(m_{\stackrel{~}{g}}^2m_{\chi _i^0}^2)^3}{2\pi m_{\stackrel{~}{g}}^3\stackrel{~}{m}^4}\left(C_{g}^{\chi _i^0}{}_{\mathrm{eff}}{}^{}\right)^2.$$
(90)
The use of the effective coefficient $`C_{g}^{\chi _i^0}{}_{\mathrm{eff}}{}^{}`$ defined in eq. (62) allows us to reproduce the complete one–loop result when the resummation is switched off.
### Decays into goldstino:
the decay width into goldstino and quarks of the first and second generation is:
$$\mathrm{\Gamma }_{\stackrel{~}{G}\overline{q}q}=\frac{m_{\stackrel{~}{g}}^5}{192\pi ^3F^2}C_{1}^{\stackrel{~}{G}}{}_{}{}^{\mathrm{\hspace{0.17em}2}},$$
(91)
where we have summed over all four light quark flavours.
The gluino decay width into goldstino and third–generation quarks is as in eq. (79), with $`\stackrel{~}{m}^4`$ replaced by $`F^2`$. The squared decay amplitude, which has to be integrated numerically on the $`(s_{13},s_{23})`$ plane, is given by
$`\overline{\left|\right|^2}`$ $`=`$ $`\left(C_{q_L}^{\stackrel{~}{G}}{}_{}{}^{2}+C_{q_R}^{\stackrel{~}{G}}{}_{}{}^{2}\right)\left[(m_{\stackrel{~}{g}}^2+m_q^2s_{13})(s_{13}m_q^2)+(m_{\stackrel{~}{g}}^2+m_q^2s_{23})(s_{23}m_q^2)\right]`$ (92)
$`+`$ $`4C_{q_L}^{\stackrel{~}{G}}C_{q_R}^{\stackrel{~}{G}}m_q^2(s_{13}+s_{23}2m_q^2)(q=t,b),`$
where
$$C_{t_L}^{\stackrel{~}{G}}=C_{b_L}^{\stackrel{~}{G}}=C_2^{\stackrel{~}{G}},C_{t_R}^{\stackrel{~}{G}}=C_3^{\stackrel{~}{G}},C_{b_R}^{\stackrel{~}{G}}=C_4^{\stackrel{~}{G}}.$$
(93)
Finally, the gluino decay width into gluon and goldstino is:
$$\mathrm{\Gamma }_{\stackrel{~}{G}g}=\frac{m_{\stackrel{~}{g}}^3}{2\pi F^2}C_{5}^{\stackrel{~}{G}}{}_{}{}^{\mathrm{\hspace{0.17em}2}}.$$
(94)
## Acknowledgements
We thank M. Toharia and J. Wells for precious help in the comparison with the results of ref. . We also thank M. Gorbahn, U. Haisch and P. Richardson for useful discussions. P. S. thanks the CERN Theory Division and INFN, Sezione di Torino for hospitality during the completion of this work. The work of P. G. is supported in part by the EU grant MERG-CT-2004-511156 and by MIUR under contract 2004021808-009.
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# Eigenvalues and degree deviation in graphs
## 1 Introduction
Our notation is standard (e.g., see , , and ); in particular, all graphs are defined on the vertex set $`\{1,2,\mathrm{},n\}=\left[n\right]`$ and $`G(n,m)`$ stands for a graph with $`n`$ vertices and $`m`$ edges. We write $`\mathrm{\Gamma }\left(u\right)`$ for the set of neighbors of the vertex $`u`$ and set $`d\left(u\right)=\left|\mathrm{\Gamma }\left(u\right)\right|.`$ Given a graph $`G`$ of order $`n,`$ we assume that the eigenvalues of the adjacency matrix of $`G`$ are ordered as $`\mu \left(G\right)=\mu _1\left(G\right)\mathrm{}\mu _n\left(G\right)`$. As usual, $`\overline{G}`$ denotes the complement of a graph $`G.`$
Collatz and Sinogowitz showed that $`\mu \left(G\right)2m/n`$ for every graph $`G=G(n,m).`$ Since equality holds if and only if $`G`$ is regular, they proposed the value $`ϵ\left(G\right)=\mu \left(G\right)2m/n`$ as a relevant measure of irregularity of $`G`$. Two other closely related measures of graph irregularity are the functions
$`var\left(G\right)`$ $`={\displaystyle \frac{1}{n}}{\displaystyle \underset{uV\left(G\right)}{}}\left(d\left(u\right){\displaystyle \frac{2m}{n}}\right)^2,`$
$`s\left(G\right)`$ $`={\displaystyle \underset{uV\left(G\right)}{}}\left|d\left(u\right){\displaystyle \frac{2m}{n}}\right|.`$
Bell compared $`ϵ\left(G\right)`$ to $`var\left(G\right)`$ and showed that none of them could be preferred to the other one as a measure of irregularity. He did not, however, give explicit inequalities between $`ϵ\left(G\right)`$ and $`var\left(G\right)`$. In this note we prove that for every graph $`G`$ with $`n`$ vertices and $`m`$ edges,
$$\frac{var\left(G\right)}{2\sqrt{2m}}\mu \left(G\right)\frac{2m}{n}\sqrt{s\left(G\right)}$$
(1)
Thus, in view of
$$\frac{s^2\left(G\right)}{n^2}var\left(G\right)s\left(G\right),$$
we also have
$$\frac{s^2\left(G\right)}{2n^2\sqrt{2m}}\mu \left(G\right)\frac{2m}{n}\sqrt{var\left(G\right)}.$$
In addition we derive similar inequalities specifically for bipartite graphs.
Another well-known inequality involving graph eigenvalues is
$$\mu _k\left(G\right)+\mu _{nk+1}\left(\overline{G}\right)1,$$
(2)
holding for every graph $`G`$ of order $`n`$ and every $`k=1,\mathrm{},n1.`$ Note that if $`G`$ is regular, equality holds in (2) but the converse is not always true (e.g., $`G=K_{a,b},`$ $`b>a>2,`$ $`k=2)`$. A natural problem is to find a lower bound on $`\mu _k\left(G\right)+\mu _{nk+1}\left(\overline{G}\right)`$ implying explicit equality in (2) for regular $`G`$. In this note we show that for every $`k=1,\mathrm{},n1,`$
$$\mu _k\left(G\right)+\mu _{nk+1}\left(\overline{G}\right)12\sqrt{2s\left(G\right)}.$$
(3)
We show that inequalities (1) and (3) are tight up to a constant factor.
Finally we prove that for every graph $`G`$ of order $`n`$,
$$\mu _n\left(G\right)+\mu _n\left(\overline{G}\right)1\frac{s^2\left(G\right)}{n^3},$$
(4)
implying that for any highly irregular graph $`G`$ either $`\mu _n\left(G\right)`$ or $`\mu _n\left(\overline{G}\right)`$ must be large in absolute value.
Let us note that these results are readily applicable to the study of quasirandom graph properties.
The rest of the note is organized as follows. In Section 2 we describe algorithms for regularizing graphs with few edge changes. Section 3 contains basic results about spectra of blown-up graphs. In Sections 4, 5, and 6 we prove inequalities (1), (3), and (4).
## 2 Efficient regularization
Consider the following natural problem: given a graph $`G,`$ what is the minimum number of edges $`\rho \left(G\right)`$ that must be changed to obtain a regular graph. Writing $`A\left(G\right)`$ for the adjacency matrix of a graph $`G,`$ we see that
$$\rho \left(G\right)=\frac{1}{2}\mathrm{min}\left\{A\left(G\right)A\left(R\right)_2:R\text{ is regular graph of order }v\left(G\right)\right\}.$$
It is almost certain that the problem of estimating $`\rho \left(G\right)`$ has been raised and solved in the literature but, lacking a proper reference, we shall solve it from scratch.
We first show that there exists a graph $`R^{}`$ whose degrees differ by at most one and such that
$$A\left(G\right)A\left(R^{}\right)_2s\left(G\right).$$
Next we find a regular graph $`R`$ such that
$$A\left(R\right)A\left(R^{}\right)_2<3n.$$
Finally we show that for every graph $`G,`$
$$\rho \left(G\right)s\left(G\right)/2$$
implying that our upper bounds on $`\rho \left(G\right)`$ are not too far from the best possible ones.
### 2.1 Rough regularization
The main result in this section is the following theorem.
###### Theorem 1
For every graph $`G=G(n,m),`$ there exists a graph $`R=G(n,m)`$ such that $`\mathrm{\Delta }\left(R\right)\delta \left(R\right)+1`$ and $`R`$ differs from $`G`$ in at most $`s\left(G\right)`$ edges. In particular, if $`2m/n`$ is integer then $`R`$ is $`\left(2m/n\right)`$-regular.
Proof We shall describe a simple algorithm that produces the graph $`R`$ by deleting and adding edges of $`G.`$ Set $`d=2m/n`$.
Step 1
*While* $`\delta \left(G\right)<d`$ *and* $`\mathrm{\Delta }\left(G\right)>d+1`$ *select* $`u,v`$ *with* $`d\left(u\right)=\delta \left(G\right)`$ *and* $`d\left(v\right)=\mathrm{\Delta }\left(G\right)`$*. Since* $`\mathrm{\Gamma }\left(v\right)\backslash \mathrm{\Gamma }\left(u\right)\mathrm{},`$ *there exists* $`w\mathrm{\Gamma }\left(v\right)\backslash \mathrm{\Gamma }\left(u\right);`$ *delete the edge* $`vw`$ *and add the edge* $`uw.`$
Write $`G^{}`$ for the graph obtained upon exiting Step 1. Since Step 1 is iterated as long as $`\delta \left(G\right)<d`$ and $`\mathrm{\Delta }\left(G\right)>d+1`$, we have either $`\delta \left(G^{}\right)=d`$ or $`\mathrm{\Delta }\left(G^{}\right)=d+1;`$ we may assume $`\delta \left(G^{}\right)=d,`$ since the other case is reduced to this one by considering $`\overline{G^{}}`$.
If $`\mathrm{\Delta }\left(G^{}\right)d+1`$ then terminate the procedure with $`R=G^{}.`$ Otherwise write $`A`$ for the set of vertices of degree $`d,`$ $`B`$ for the set of vertices of degree $`d+1,`$ and $`C`$ for the set of vertices of degree $`d+2`$ or higher.
Step 2
*While* $`C\mathrm{}`$ *select* $`uA,`$$`vC.`$ *Since* $`\left|\mathrm{\Gamma }\left(v\right)\right|>\left|\mathrm{\Gamma }\left(u\right)\right|,`$ *we may select* $`w\mathrm{\Gamma }\left(v\right)\backslash \mathrm{\Gamma }\left(u\right);`$ *delete the edge* $`vw`$ *and add the edge* $`uw.`$
Write $`R`$ for the graph obtained after executing Step 2. Let $`G^{},A,B,C`$ be as defined prior to Step 2; set $`\left|A\right|=k,`$ $`\left|C\right|=s`$. Each iteration in Step 1 changes two edges and decreases $`s\left(G\right)`$ by 2; therefore, after the execution of Step 1, at most $`s\left(G\right)s\left(G^{}\right)`$ edges of $`G`$ are changed. Set
$$l=\underset{uC}{}\left(d\left(u\right)d1\right).$$
Each iteration in Step 2 changes two edges and decreases $`l`$ by 1; therefore, there are $`l`$ iterations in Step 2 and at most $`2l`$ edges are changed. To complete the proof we have to show that $`S\left(G^{}\right)2l.`$ From
$`2m/n`$ $`={\displaystyle \frac{1}{n}}{\displaystyle \underset{uV\left(G\right)}{}}d\left(u\right)={\displaystyle \frac{1}{n}}{\displaystyle \underset{uA}{}}d\left(u\right)+{\displaystyle \frac{1}{n}}{\displaystyle \underset{uB}{}}d\left(u\right)+{\displaystyle \frac{1}{n}}{\displaystyle \underset{uC}{}}d\left(u\right)`$
$`={\displaystyle \frac{kd+\left(nks\right)\left(d+1\right)+s\left(d+1\right)+l}{n}}=d+{\displaystyle \frac{nk+l}{n}}`$
it follows that $`k>l.`$ Furthermore,
$`s\left(G^{}\right)2l`$ $`={\displaystyle \underset{uV}{}}\left|d_G^{}\left(u\right)2m/n\right|2l`$
$`={\displaystyle \underset{uA}{}}\left|d_G^{}\left(u\right)2m/n\right|+{\displaystyle \underset{uB}{}}\left|d_G^{}\left(u\right)2m/n\right|+{\displaystyle \underset{uC}{}}\left|d_G^{}\left(u\right)2m/n\right|2l`$
$`=k{\displaystyle \frac{nk+l}{n}}+\left(nks\right){\displaystyle \frac{kl}{n}}+s{\displaystyle \frac{kl}{n}}l=2{\displaystyle \frac{\left(kl\right)\left(nk\right)}{n}}>0,`$
completing the proof. $`\mathrm{}`$
#### 2.1.1 Rough regularization of bipartite graphs
Call a bipartite graph *semiregular* if vertices belonging to the same vertex class have equal degrees.
Let $`G`$ be a bipartite graph and $`A,B`$ be its vertex classes, $`\left|A\right|=a,`$ $`\left|B\right|=b.`$ Define the function
$$s_2\left(G\right)=\underset{uA}{}\left|d\left(u\right)\frac{m}{a}\right|+\underset{uB}{}\left|d\left(u\right)\frac{m}{b}\right|;$$
$`s_2\left(G\right)`$ is the equivalent to $`s\left(G\right)`$ for bipartite graphs. Clearly, $`s_2\left(G\right)=0`$ if and only if $`G`$ is semiregular.
Modifying slightly the proof of Theorem 1 we obtain the following special case for bipartite graphs.
###### Theorem 2
For every bipartite graph $`G=G(n,m)`$ with vertex classes $`A,B,`$ there exists a bipartite graph $`R=G(n,m)`$ with the same vertex classes such that:
*(i)* $`\left|d_R\left(u\right)d_R\left(v\right)\right|1`$ for every $`u,v`$ belonging to the same vertex class;
*(ii)* $`R`$ differs from $`G`$ in at most $`s_2\left(G\right)`$ edges.
In particular, if $`m/\left|A\right|`$ and $`m/\left|B\right|`$ are integer then $`R`$ is semiregular.
### 2.2 Fine regularization
If we allow $`m`$ to change, we may further regularize the graph $`R`$ obtained in Theorem 1.
###### Theorem 3
Let the degrees of a graph $`G=G(n,m)`$ be either $`d`$ or $`d+1.`$ There exists an $`r`$-regular graph $`R`$ such that either $`r=d`$ or $`r=d+1,`$ and $`R`$ differs from $`G`$ in at most $`3n/2`$ edges.
Proof Write $`A`$ for the set of vertices of degree $`d+1`$ and $`B`$ for $`V\left(G\right)\backslash A.`$ Clearly either $`\left|A\right|`$ or $`\left|B\right|`$ is even. We shall assume that $`\left|A\right|`$ is even, otherwise we may apply the argument to the complementary graph. Set $`a=\left|A\right|.`$ Our goal is to construct a $`d`$-regular graph by changing at most $`3a/2`$ edges. We shall describe a procedure constructing $`R.`$
Step 1
*While* $`E\left(A\right)\mathrm{},`$ *select* $`uvE\left(A\right)`$ *and remove it.*
Step 2.
*While* $`A\mathrm{},`$ *select two distinct* $`u,vA`$ *and two disjoint vertices* $`t\mathrm{\Gamma }\left(v\right),`$$`w\mathrm{\Gamma }\left(u\right).`$ *Delete the edges* $`uw`$ *and* $`vt;`$ *add the edge* $`wt.`$
The iteration in Step 2 may always be executed since, for every two distinct $`u,vA,`$ there exist disjoint vertices $`t\mathrm{\Gamma }\left(v\right)`$ and $`w\mathrm{\Gamma }\left(u\right).`$ Indeed, if $`\mathrm{\Gamma }\left(u\right)\mathrm{\Gamma }\left(v\right),`$ select $`w\mathrm{\Gamma }\left(u\right)\backslash \mathrm{\Gamma }\left(v\right).`$ Since $`d\left(w\right)=d<\left|\mathrm{\Gamma }\left(v\right)\right|,`$ there exists $`t\mathrm{\Gamma }\left(v\right)`$ that is disjoint from $`w`$ and the assertion is proved. If $`\mathrm{\Gamma }\left(u\right)=\mathrm{\Gamma }\left(v\right)`$ then $`\mathrm{\Gamma }\left(u\right)`$ cannot induce a complete graph, since $`\mathrm{\Gamma }\left(u\right)B`$ and so the vertices in $`\mathrm{\Gamma }\left(u\right)`$ have degree $`d.`$
Each iteration in Step 1 removes two vertices from $`A`$ and changes two edges. Each iteration in Step 2 removes two vertices from $`A`$ and changes three edges. Therefore, after changing at most $`3\left|A\right|/2`$ edges, we obtain a $`d`$-regular graph $`R`$, as claimed. $`\mathrm{}`$
### 2.3 Optimal regularization
Summarizing Theorems 1 and 3, we obtain the following corollary.
###### Corollary 4
For every graph $`G`$ of order $`n,`$
$$\rho \left(G\right)s\left(G\right)+3n/2.$$
It turns out that this bound is quite close to the optimal one, no matter what the graph $`G`$ is. We shall show that
$$\rho \left(G\right)s\left(G\right)/2.$$
Let $`R`$ be $`r`$-regular graph with $`V\left(R\right)=V\left(G\right).`$ For every vertex $`vV\left(G\right),`$ we have
$$\left|\left(\mathrm{\Gamma }_G\left(u\right)\backslash \mathrm{\Gamma }_R\left(u\right)\right)\left(\mathrm{\Gamma }_R\left(u\right)\backslash \mathrm{\Gamma }_G\left(u\right)\right)\right|d\left(u\right)+r2\mathrm{min}(d,r)\left|d\left(u\right)r\right|.$$
Hence, summing over all vertices $`vV\left(G\right)`$ we find that
$$2\rho \left(G\right)A\left(G\right)A\left(R\right)_2\left|d\left(u\right)r\right|s\left(G\right),$$
as claimed.
We note without proof that $`\rho \left(K_{a,b}\right)3s\left(K_{a,b}\right)/4`$.
## 3 The spectra of blown-up graphs
In this section we introduce two operations on graphs and consider how they affect graph spectra.
Let $`G=G(n,m)`$ and $`t>0`$ be integer. Write $`G^{\left(t\right)}`$ for the graph obtained by replacing each vertex $`uV\left(G\right)`$ by a set $`V_u`$ of $`t`$ vertices and joining $`xV_u`$ to $`yV_v`$ if and only if $`uvE\left(G\right).`$ Notice that $`v\left(G^{\left(t\right)}\right)=tn.`$ The following theorem holds.
###### Theorem 5
The eigenvalues of $`G^{\left(t\right)}`$ are $`t\mu _1\left(G\right),\mathrm{},t\mu _n\left(G\right)`$ together with $`n\left(t1\right)`$ additional $`0`$’s.
Set $`G^{\left[t\right]}=\overline{\overline{G}^{\left(t\right)}},`$ i.e., $`G^{\left[t\right]}`$ is obtained from $`G^{\left(t\right)}`$ by joining all vertices within $`V_u`$ for every $`uV\left(G\right);`$ note also that $`\overline{G^{\left(t\right)}}=\overline{G}^{\left[t\right]}.`$ The following theorem holds.
###### Theorem 6
The eigenvalues of $`G^{\left[t\right]}`$ are $`t\mu _1\left(G\right)+t1,\mathrm{},t\mu _n\left(G\right)+t1`$ together with $`n\left(t1\right)`$ additional $`\left(1\right)`$’s.
## 4 Bounds on $`\mu \left(G\right)`$
In this section we shall prove inequalities (1). Recall first the inequality
$$\mu ^2\left(G\right)\frac{1}{n}\underset{uV\left(G\right)}{}d^2\left(u\right),$$
(5)
due to Hofmeister and observe that Stanley’s inequality
$$\mu \left(G\right)1/2+\sqrt{2m+1/4}$$
implies
$$\mu ^2\left(G\right)2m.$$
(6)
We thus find that
$`2\sqrt{2m}\left(\mu \left(G\right)2m/n\right)`$ $`2\mu \left(G\right)\left(\mu \left(G\right)2m/n\right)\mu ^2\left(G\right)\left(2m/n\right)^2`$
$`{\displaystyle \frac{1}{n}}{\displaystyle \underset{uV\left(G\right)}{}}d^2\left(u\right)\left(2m/n\right)^2=var\left(G\right),`$
obtaining the lower bound in (1). To prove the upper bound we need the following proposition.
###### Proposition 7
If $`G_1`$ and $`G_2`$ are graphs with $`V\left(G_1\right)=V\left(G_2\right)`$ then
$$\mu \left(G_1\right)\mu \left(G_2\right)\sqrt{2\left|E\left(G_1\right)\backslash E\left(G_2\right)\right|}.$$
Proof Setting $`G^{}=(V\left(G_1\right),E\left(G_1\right)E\left(G_2\right)),`$ $`G^{\prime \prime }=(V\left(G_1\right),E\left(G_1\right)\backslash E\left(G_2\right)),`$ from Weyl’s inequalities (, p. 181), we have
$$\mu \left(G_1\right)\mu \left(G^{}\right)\mu \left(G_2\right)+\mu \left(G^{\prime \prime }\right).$$
By (6), we have,
$$\mu \left(G^{\prime \prime }\right)\sqrt{2\left|E\left(G_1\right)\backslash E\left(G_2\right)\right|},$$
completing the proof. $`\mathrm{}`$
We shall deduce the upper bound in (1) essentially from Theorem 1.
###### Theorem 8
For every graph $`G=G(n,m),`$
$$\mu \left(G\right)2m/n\sqrt{s\left(G\right)}.$$
Proof Theorem 1 implies that there exists a graph $`R=G(n,m)`$ such that $`\mathrm{\Delta }\left(R\right)\delta \left(r\right)+1`$ and $`R`$ differs from $`G`$ in at most $`s\left(G\right)`$ edges. Since $`e\left(R\right)=e\left(G\right)`$ it follows that $`\left|E\left(G\right)\backslash E\left(R\right)\right|=\left|E\left(R\right)\backslash E\left(G\right)\right|`$ and so $`2\left|E\left(G\right)\backslash E\left(R\right)\right|s\left(G\right).`$ Hence, by Proposition 7,
$$\mu \left(G\right)2m/n\mu \left(G\right)2m/n+1\mu \left(G\right)\mu \left(R\right)+11+\sqrt{s\left(G\right)}.$$
(7)
Notice that $`v\left(G^{\left(t\right)}\right)=tn,`$ $`e\left(G^{\left(t\right)}\right)=t^2m,`$ and $`s\left(G^{\left(t\right)}\right)=t^2s\left(G\right).`$ Applying Theorem 5, we also see that
$$\mu \left(G^{\left(t\right)}\right)=t\mu \left(G\right).$$
From (7) it follows that
$$\left(\mu \left(G\right)2m/n\right)t=\mu \left(G^{\left(t\right)}\right)2e\left(G^{\left(t\right)}\right)/v\left(G^{\left(t\right)}\right)1+\sqrt{s\left(G^{\left(t\right)}\right)}=1+t\sqrt{s\left(G\right)}.$$
Hence, dividing by $`t`$ and letting $`t`$ tend to infinity, the desired inequality follows. $`\mathrm{}`$
### 4.1 Tightness of inequalities (1)
It is natural to ask how large $`c`$ could be so that the inequality
$$\mu \left(G\right)\frac{2m}{n}c\frac{s^2\left(G\right)}{n^2\sqrt{m}}$$
holds for every graph $`G=G(n,m).`$ Taking the graph $`G=K_{n,n+1}`$ for $`n`$ large enough, we see that $`c`$ may be at most $`1/2`$.
Similarly, let $`c`$ be such that the inequality
$$\mu \left(G\right)2m/nc\sqrt{s\left(G\right)}$$
holds for every graph $`G=G(n,m).`$ Taking $`G=K_nK_1`$ we see that $`c`$ must be at least $`1/\sqrt{2}.`$
We venture the following conjecture.
###### Conjecture 9
For every graph $`G`$ of sufficiently large order $`n`$ and size $`m`$,,
$$\frac{s^2\left(G\right)}{2n^2\sqrt{m}}\mu \left(G\right)\frac{2m}{n}\sqrt{s\left(G\right)/2}.$$
### 4.2 Bounds on $`\mu \left(G\right)`$ when $`G`$ is bipartite
It is possible to modify inequalities (1) to better suit bipartite graphs.
Let $`G`$ be a bipartite graph and $`A,B`$ be its vertex classes, $`\left|A\right|=a,`$ $`\left|B\right|=b.`$ Then, by Rayleigh’s principle we have,
$$\mu \left(G\right)e\left(G\right)/\sqrt{ab}.$$
A careful analysis shows that equality is possible if and only if $`G`$ is semiregular. In fact the following theorem holds.
###### Theorem 10
For every bipartite graph $`G`$ with vertex classes $`A,B,`$
$$\frac{s_2^2\left(G\right)}{2n^2\sqrt{\left|A\right|\left|B\right|}}\mu \left(G\right)\frac{e\left(G\right)}{\sqrt{\left|A\right|\left|B\right|}}\sqrt{\frac{s_2\left(G\right)}{2}}.$$
Proof Let $`\left|A\right|=a,`$ $`\left|B\right|=b,`$ $`e\left(G\right)=m,`$ $`v\left(G\right)=n.`$ We start with the proof of the first inequality. By the AM-QM inequality we have
$`{\displaystyle \underset{uA}{}}\left|d\left(u\right){\displaystyle \frac{m}{a}}\right|`$ $`\sqrt{a{\displaystyle \underset{uA}{}}\left(d\left(u\right){\displaystyle \frac{m}{a}}\right)^2},`$
$`{\displaystyle \underset{uB}{}}\left|d\left(u\right){\displaystyle \frac{m}{b}}\right|`$ $`\sqrt{b{\displaystyle \underset{uB}{}}\left(d\left(u\right){\displaystyle \frac{m}{b}}\right)^2}.`$
Hence, by Cauchy-Schwarz and inequality (5), we find that,
$`s_2\left(G\right)`$ $`\sqrt{n}\sqrt{{\displaystyle \underset{uA}{}}\left(d\left(u\right){\displaystyle \frac{m}{a}}\right)^2+{\displaystyle \underset{uB}{}}\left(d\left(u\right){\displaystyle \frac{m}{b}}\right)^2}=\sqrt{n}\sqrt{{\displaystyle \underset{uV\left(G\right)}{}}\left(d^2\left(u\right){\displaystyle \frac{m^2n}{ab}}\right)}`$
$`n\sqrt{\mu ^2\left(G\right){\displaystyle \frac{m^2}{ab}}}n\sqrt{\left(\mu \left(G\right){\displaystyle \frac{m}{\sqrt{ab}}}\right)\left(2\sqrt{ab}\right)},`$
proving the first inequality.
To prove the second inequality we first note the equivalent of Proposition 7 for bipartite graphs: if $`G_1`$ and $`G_2`$ are bipartite graphs with the same vertex classes then
$$\mu \left(G_1\right)\mu \left(G_2\right)\sqrt{\left|E\left(G_1\right)\backslash E\left(G_2\right)\right|}.$$
Note that the coefficient $`2`$ under the square root is missing here, since $`\mu \left(G\right)\sqrt{e\left(G\right)}`$ for bipartite $`G`$ (Cvetković , also , p. 92 Theorem 3.19).
Theorem 2 implies that there exists a graph $`R=G(n,m)`$ with vertex classes $`A,B`$ such that $`\left|d_R\left(u\right)d_R\left(v\right)\right|1`$ for every $`u,v`$ belonging to the same vertex class and $`R`$ differs from $`G`$ in at most $`s_2\left(G\right)`$ edges. Since $`e\left(R\right)=e\left(G\right)`$ it follows that $`\left|E\left(G\right)\backslash E\left(R\right)\right|=\left|E\left(R\right)\backslash E\left(G\right)\right|`$ and so $`2\left|E\left(G\right)\backslash E\left(R\right)\right|s_2\left(G\right).`$ Hence, by Proposition 7,
$$\mu \left(G\right)\mu \left(R\right)\sqrt{\frac{s_2\left(G\right)}{2}}.$$
Applying the inequality $`\mu \left(G\right)\mathrm{max}_{uvE\left(G\right)}\sqrt{d\left(u\right)d\left(v\right)},`$ due to Berman and Zhang , we find that
$$\mu \left(R\right)\sqrt{\left(\frac{m}{a}+1\right)\left(\frac{m}{b}+1\right)}\sqrt{\frac{m^2}{ab}+\frac{mn}{ab}+1}<\frac{m}{\sqrt{ab}}+\sqrt{n+1}$$
and so,
$$\mu \left(G\right)\frac{m}{\sqrt{ab}}\sqrt{\frac{s_2\left(G\right)}{2}}+\sqrt{n+1}.$$
Now, applying the final argument from the proof of Theorem 8, the desired inequality follows. $`\mathrm{}`$
## 5 A lower bound on $`\mu _k\left(G\right)+\mu _{nk+1}\left(\overline{G}\right)`$
The main goal of this section is the proof of inequality (2). By Weyl’s inequalities (, p. 181), for every graph $`G`$ of order $`n,`$ we have
$$\mu _k\left(G\right)+\mu _{nk+1}\left(\overline{G}\right)\mu _k\left(K_n\right)=1.$$
###### Theorem 11
For every $`k=1,\mathrm{},n1`$
$$\mu _k\left(G\right)+\mu _{nk+1}\left(\overline{G}\right)12\sqrt{2s\left(G\right)}$$
Proof By Corollary 4 there exists a regular graph $`R`$ that differs from $`G`$ in at most $`s\left(G\right)+3n/2`$ edges. Then, by Weyl’s inequalities,
$`\mu _k\left(A\left(G\right)\right)+\mu _1\left(A\left(R\right)A\left(G\right)\right)`$ $`\mu _k\left(A\left(R\right)\right),`$
$`\mu _{nk+1}\left(A\left(\overline{G}\right)\right)+\mu _1\left(A\left(\overline{R}\right)A\left(\overline{G}\right)\right)`$ $`\mu _{nk+1}\left(A\left(\overline{R}\right)\right).`$
Furthermore, by
$`\mu _1\left(A\left(R\right)A\left(G\right)\right)`$ $`\sqrt{A\left(R\right)A\left(G\right)_2}=\sqrt{2s\left(G\right)+3n}`$
$`\mu _1\left(A\left(\overline{R}\right)A\left(\overline{G}\right)\right)`$ $`\sqrt{A\left(\overline{R}\right)A\left(\overline{G}\right)_2}=\sqrt{2s\left(G\right)+3n},`$
we find that
$`\mu _k\left(G\right)+\mu _{nk+1}\left(\overline{G}\right)`$ $`\mu _k\left(A\left(R\right)\right)+\mu _{nk+1}\left(A\left(\overline{R}\right)\right)2\sqrt{2s\left(G\right)+3n}`$
$`=12\sqrt{2s\left(G\right)+3n}.`$
Suppose now that $`t`$ is sufficiently large and consider the graphs $`G^{\left(t\right)}`$ and $`\overline{G^{\left(t\right)}}.`$ By Theorem 5 we have
$$\mu _k\left(G^{\left(t\right)}\right)=t\mu _k\left(G\right).$$
Similarly in view of and $`\overline{G^{\left(t\right)}}=\overline{G}^{\left[t\right]}`$ and Theorem 6,
$$\mu _{ntk+1}\left(\overline{G^{\left(t\right)}}\right)\mathrm{min}\{t\mu _{nk+1}\left(\overline{G}\right)+t1,1\}$$
Since, $`s\left(G^{\left(t\right)}\right)=t^2s\left(G\right),`$ we see that
$`t\mu _k\left(G\right)+t\mu _{nk+1}\left(\overline{G}\right)`$ $`\mu _k\left(G^{\left(t\right)}\right)+\mu _k\left(\overline{G^{\left(t\right)}}\right)t+1t2\sqrt{2s\left(G^{\left(t\right)}\right)+3nt}`$
$`=t2t\sqrt{2s\left(G\right)+3n/t}.`$
Dividing by $`t`$ and letting $`t`$ tend to infinity, we obtain the desired inequality. $`\mathrm{}`$
For the graph $`G=K_{1,n}`$ we have $`s\left(G\right)=2\frac{n1}{n+1}`$ and $`\mu _{n+1}\left(G\right)+\mu _2\left(\overline{G}\right)=1\sqrt{n}.`$ Hence,
$$\mu _{n+1}\left(G\right)+\mu _2\left(\overline{G}\right)=1\left(\frac{1}{\sqrt{2}}+o\left(1\right)\right)\sqrt{s\left(G\right)},$$
implying that inequality (3) is tight up to a constant factor less than $`4.`$
## 6 An upper bound on $`\mu _n\left(G\right)+\mu _n\left(\overline{G}\right)`$
The main result in this section is the proof of inequality (4). We start with an auxiliary result.
###### Lemma 12
For every graph $`G`$ of order $`n`$ there exists an $`n/2`$-set $`SV\left(G\right)`$ such that
$$e\left(V\left(G\right)\backslash S\right)e\left(S\right)\frac{1}{2}s\left(G\right).$$
Proof Note first that for any $`a`$ we have
$$\underset{i=1}{\overset{n}{}}\left|d_ia\right|\underset{i=1}{\overset{n}{}}\left|d_i\frac{2m}{n}\right|=s\left(G\right).$$
Let $`d\left(1\right)d\left(2\right)\mathrm{}.d\left(n\right)`$ be the degree sequence of $`G`$ and set $`V=\left[n\right].`$ For every $`1kn,`$ letting $`S=[k],`$ we have
$$\underset{uV\backslash S}{}d\left(u\right)\underset{uS}{}d\left(u\right)=2e\left(S\right)+e(S,V\backslash S)2e\left(V\backslash S\right)e(S,V\backslash S)=2e\left(V\backslash S\right)2e\left(S\right).$$
Assume first $`n`$ even, $`n=2k`$. Letting $`a=\left(d\left(k\right)+d\left(k+1\right)\right)/2`$ and $`S=[k],`$ we have
$$\underset{uV\backslash S}{}d\left(u\right)\underset{uS}{}d\left(u\right)=\underset{uV\backslash S}{}\left(d\left(u\right)a\right)+\underset{uS}{}\left(ad\left(u\right)\right)=\underset{uV}{}\left|d_ia\right|s\left(G\right),$$
proving the assertion for even $`n.`$
Let now $`n`$ be odd, $`n=2k+1`$. Letting $`a=d_{k+1}`$ and $`S=[k],`$ we have
$$\underset{uV\backslash S}{}d\left(u\right)\underset{uS}{}d\left(u\right)=\underset{uV\backslash S}{}\left(d\left(u\right)a\right)+\underset{uS}{}\left(ad\left(u\right)\right)=\underset{uV}{}\left|d\left(u\right)a\right|s\left(G\right),$$
proving the assertion for odd $`n`$ as well. $`\mathrm{}`$
###### Theorem 13
For every graph $`G`$ of order $`n,`$
$$\mu _n\left(G\right)+\mu _n\left(\overline{G}\right)1\frac{s^2\left(G\right)}{n^3}.$$
Proof From the interlacing theorem of Haemers (see, e.g., , ), for every bipartition of $`V\left(G\right)=V_1V_2`$ we have
$$\mu _n\left(G\right)\frac{e\left(V_1\right)}{\left|V_1\right|}+\frac{e\left(V_2\right)}{\left|V_2\right|}\sqrt{\left(\frac{e\left(V_1\right)}{\left|V_1\right|}\frac{e\left(V_2\right)}{\left|V_2\right|}\right)^2+\frac{e(V_1,V_2)^2}{\left|V_1\right|\left|V_2\right|}}.$$
(8)
Assume $`n`$ even and let $`V\left(G\right)=V_1V_2`$ be a bipartition such that $`\left|V_1\right|=\left|V_2\right|=n/2,`$ and $`e\left(V_1\right)e\left(V_2\right)s\left(G\right)/2.`$ Letting $`e_1=e\left(V_1\right)`$, $`e_2=e\left(V_2\right),`$ $`e_3=e(V_1,V_2),`$ $`s=s\left(G\right),`$ from (8), after some simple algebra, we obtain
$$\frac{n}{2}\mu _n\left(G\right)e_1+e_2\sqrt{\left(e_1e_2\right)^2+e_3^2}e_1+e_2\sqrt{\frac{s^2}{4}+e_3^2}.$$
(9)
Note that $`s\left(G\right)<n^2`$ and $`e(V_1,V_2)n^2/4;`$ thus, we have
$$\frac{s^4}{9n^4}+\frac{2e_3s^2}{3n^2}+e_3^2s^2\left(\frac{1}{9}+\frac{1}{6}\right)+e_3^2\frac{s^2}{4}+e_3^2,\text{ }$$
and so,
$$\sqrt{\frac{s^2}{4}+e_3^2}\frac{s^2}{3n^2}+e_3.$$
Hence, from (9), it follows that
$$\frac{n}{2}\mu _n\left(G\right)e_1+e_2e_3\frac{s^2}{3n^2}.$$
Since $`s\left(G\right)=s\left(\overline{G}\right),`$ we see also that
$$\frac{n}{2}\mu _n\left(\overline{G}\right)\left(\genfrac{}{}{0pt}{}{n/2}{2}\right)e_1+\left(\genfrac{}{}{0pt}{}{n/2}{2}\right)e_2\frac{n^2}{4}+e_3\frac{s^2}{3n^2}$$
and hence,
$$\frac{n}{2}\left(\mu _n\left(G\right)+\mu _n\left(\overline{G}\right)\right)2\left(\genfrac{}{}{0pt}{}{n/2}{2}\right)\frac{n^2}{4}\frac{2s^2}{3n^2}=\frac{n}{2}\frac{2s^2}{3n^2},$$
proving the assertion for even $`n.`$
To prove the assertion for odd $`n`$ observe that if $`t`$ is even, for the graph $`G^{\left(t\right)}`$ we have
$$t\mu _n\left(G\right)+t\mu _n\left(\overline{G}\right)+t1=\mu _{tn}\left(G^{\left(t\right)}\right)+\mu _{tn}\left(\overline{G^{\left(t\right)}}\right)1\frac{t^4s^2}{t^3n^3}.$$
Dividing by $`t`$ and letting $`t`$ tend to infinity, the assertion follows for odd $`n`$ as well. $`\mathrm{}`$
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# Self-stabilization in certain infinite-dimensional matrix algebras
## 1. Introduction
Learning $`K`$-theory, one likely encounters stabilization of matrices, linearization of cyclic loops, and the contractibility of the pointed stable Toeplitz algebra. Stabilization of matrices is a fundamental feature of $`K`$-theory; linearization of cyclic loops is an important method to prove complex Bott periodicity; the Toeplitz algebra can also be used for the same purpose, but it is also a tool to construct classifying spaces. Although considered simple, these basic constructions are often treated in quite awkward manners. The purpose of this paper is to show that these topics can be discussed in a unified and simple way. Our statements are formulated primarily in the setting of locally convex algebras. This is not just for the sake of extreme generality but to demonstrate that concrete formulas and maps can be very successful, without using approximations. The main statements of this paper are as follows:
###### Statement 1.1 (Self-stabilization).
Assume that $`𝔄=𝒦_{}(𝔖)`$, i. e. the locally convex algebra of rapidly decreasing $`\times `$ matrices over an other locally convex algebra $`𝔖`$. Let $`r:𝒦_{}(𝔄)𝔄`$ be an isomorphism which comes from relabeling $`𝒦_\times (𝔖)`$ into $`𝒦_{}(𝔖)`$. Then there is a smooth homotopy
$$E:𝒦_{}(𝔄)\times [0,\pi /2]𝒦_{}(𝔄)$$
such that it yields a family of endomorphisms of $`𝒦_{}(𝔄)`$, which are isomorphisms for $`\theta [0,\pi /2)`$, and a closed injective endomorphism for $`\theta =\pi /2`$, with
$$E(A,0)=A,E(A,\pi /2)=\mathrm{Diag}(\mathrm{},0,0|r(A),0,0,\mathrm{});$$
cf. (1) for the diagonal notation. This statement extends to unit groups, showing that $`𝚄(𝒦_{}(𝔄))`$ can be pushed down by a homotopy into $`𝚄(𝔄𝐞_{00})`$.
(The continuous map $`\varphi :𝔄\times [0,\pi /2]_\theta 𝔅`$ is smooth in the variable $`\theta `$ if the higher partial derivatives $`_\theta ^n\varphi :𝔄\times [0,\pi /2]_\theta 𝔅`$ are still continuous functions.)
Let $`𝔄[𝗓^1,𝗓]`$ be the algebra of formal Laurent series with rapidly decreasing coefficients.
###### Statement 1.2 (Linearization of cyclic loops).
There is a smooth homotopy
$$K:𝚄(𝔄[𝗓^1,𝗓])\times [0,\pi /2]𝚄(𝒦_{}(𝔄)[𝗓^1,𝗓]),$$
such that
$$K(a(𝗓),\pi /2)=\mathrm{Diag}(\mathrm{}1,1|a(𝗓)a(1)^1,1,1,\mathrm{}),$$
but
$$K(a(𝗓),0)=𝖴(a)\mathrm{\Lambda }(𝗓,𝖰)𝖴(a)^1\mathrm{\Lambda }(𝗓,𝖰)^1;$$
where $`\mathrm{\Lambda }(𝗓,𝖰)=\mathrm{Diag}(\mathrm{},𝗓,𝗓|\mathrm{\hspace{0.17em}1},1,\mathrm{})`$ is the linear loop generated by the Hilbert transform, cf. (1), and $`𝖴(a)`$ is the matrix of multiplication by $`a(𝗓)`$, cf. (5).
###### Statement 1.3 (Toeplitz contractibility).
Let $`𝔄=𝒦_{}(𝔖)`$. Then the unit group of the pointed Toeplitz algebra over $`𝔄`$, i. e. $`𝚄(𝒯_{}(𝔄)^{\mathrm{po}})`$, is contractible.
These statements were formulated in the smooth category. However, it is often useful to work in slightly different categories. One case is when $`𝔄`$ is a $``$-algebra. In those cases, instead of the general unit group $`𝚄(𝔄)`$ of invertible elements, one should work with the group $`𝚄^{}(𝔄)`$ of unitary elements. Another type of restriction occurs in the finite perturbation category, when the algebra $`𝒦_{}(𝔖)`$ of rapidly decreasing matrices is replaced by the algebra $`𝒦_{}^\mathrm{f}(𝔖)`$ of matrices with finitely many nonzero entries, and the algebra $`𝔄[𝗓^1,𝗓]`$ of rapidly decreasing Laurent series is replaced by the algebra $`𝔄[𝗓^1,𝗓]^\mathrm{f}`$ of finite Laurent series. (Here one should be careful, because for smooth loops being finite and invertible does not generally imply that the inverse is finite.)
###### Statement 1.4.
Statements 1.11.3 restrict to the $``$-algebra and/or finite perturbation categories.
The setting of finite perturbations, may, however, be too restrictive. Let us call an element $`a(𝗓)𝚄(𝔄[𝗓^1,𝗓])`$ algebraically finite if $`a=a_s\mathrm{}a_1`$, where for each $`s`$ either $`a_s`$ or $`(a_s)^1`$ has finite Laurent series form. The algebraically finite elements of $`𝚄(𝔄[𝗓^1,𝗓])`$ fall into various finiteness classes $`F`$ depending on the length of the elements $`a_s`$ or $`(a_s)^1`$. Let $`𝔄_F`$ be the set of decompositions $`\{a_j\}_{1js}`$ compatible with $`F`$. Then Statement 1.2 can be augmented as follows:
###### Statement 1.5.
For any finiteness class $`F`$, there is a smooth homotopy
$$K_F^\mathrm{e}:𝔄_F\times [0,1]\times [0,\pi /2]𝚄(𝒦_{}(𝔄)[𝗓^1,𝗓]),$$
such that
(i) $`K_F^\mathrm{e}(\stackrel{~}{a},0,\theta )=K(a,\theta )`$;
(ii) $`K_F^\mathrm{e}(\stackrel{~}{a},1,\theta )`$ differs from $`1_{}`$ in finitely many places (depending on $`F`$);
(iii) $`K_F^\mathrm{e}(\stackrel{~}{a},h,\pi /2)`$ is constant in $`h`$;
(iv) $`K_F^\mathrm{e}(\stackrel{~}{a},h,0)=𝖴_F(\stackrel{~}{a},h)\mathrm{\Lambda }(𝗓,𝖰)𝖴_F(\stackrel{~}{a},h)^1\mathrm{\Lambda }(𝗓,𝖰)^1.`$
Here $`𝖴_F(\stackrel{~}{a},h)`$ differs from $`𝖴(a)`$ in a rapidly decreasing matrix.
In particular, $`K_F^\mathrm{e}(\stackrel{~}{a},1,0)`$ yields a finite linearization of $`a(𝗓)a(1)^1`$.
These statements are known, but in lesser generality, in various ways: Statement 1.1, as stated here in the smooth category (however, see 1.4), follows from Cuntz, , Section 2. Statement 1.2 is a quantitative version of the well-known linearization technique of Atiyah and Bott, ; but much resembling to the formulas of Pressley and Segal, , Ch. 6, who work with Hilbert-Schmidt matrices, instead of rapidly decreasing ones. Statement 1.3 comes from the original Toeplitz argument of Cuntz, , originally stated in the context of $`C^{}`$-algebras, but subsequently adapted to the smooth case, cf. also . One can also find some explicit homotopies in . Statement 1.4 is useful, because $``$-algebras are prominent in operator algebraic discussions; and the finite perturbation category is the technically easiest setting to provide large contractible spaces for the purposes of algebraic topology. Statement 1.5 amounts to an explicit computation in the less functorial but more concrete setting of .
The constructions presented here are improved versions of some constructions which can be found in the author’s thesis . The author indebted to Prof. Richard B. Melrose, his advisor, for helpful discussions. In fact, much of this content was motivated by the geometric idea of Melrose, Rochon . The author would also like to thank Prof. Joachim Cuntz, who called his attention to some related papers, and Prof. Balázs Csikós, for some useful advices.
## 2. A general framework for computations
If $`𝔄`$ is a not necessarily unital algebra, then one can consider the semigroup $`1+𝔄`$, with elements of form $`1+a`$, $`(a𝔄)`$, which multiply as $`(1+a)(1+b)=1+(a+b+ab)`$. If $`𝔄`$ is unital, then it is customary to identify $`𝔄`$ and $`1+𝔄`$ by the recipe $`a𝔄\mathrm{\hspace{0.17em}1}(1_𝔄a)1+𝔄`$. This is also the situation if there is a natural identity element which can be associated to $`𝔄`$, like the identity matrix in the case of matrix algebras. The unit group $`𝚄(𝔄)`$ of $`𝔄`$ is the unit group of the semigroup $`1+𝔄`$, i. e., it is the group of pairs $`(1+a,1+b)(1+𝔄)\times (1+𝔄)`$ such that $`(1+a)(1+b)=(1+b)(1+a)=1`$; they multiply $`(1+a_1,1+b_1)(1+a_2,1+b_2)=((1+a_1)(1+a_2),(1+b_2)(1+b_1))`$. If $`𝔄`$ is a topological ring, then the natural topology on $`𝚄(𝔄)`$ comes from the product topology of $`(1+𝔄)\times (1+𝔄)`$ by restriction. As $`1+a`$ determines $`1+b`$, we write “$`1+a`$” instead of “$`(1+a,1+b)`$”. If $`\varphi :𝔄𝔅`$ is a homomorphism, then it induces a homomorphism $`𝚄\varphi :𝚄(𝔄)𝚄(𝔅)`$ defined by $`1+a1+\varphi (a)`$. We will write $`\varphi `$ instead of $`𝚄\varphi `$.
In what follows, a “locally convex vector space $`𝔄`$” means a sequentially complete, Hausdorff, locally convex vector space $`𝔄`$. The completeness is essential for analytic purposes. If the topology of $`𝔄`$ is induced by a set $`\mathrm{\Pi }_𝔄`$ of seminorms, then we assume that any positive integral combination of these seminorms also belongs to the generating seminorm set. A locally convex algebra $`𝔄`$ is a locally convex vector space with continuous bilinear multiplication. So, for each seminorm $`p\mathrm{\Pi }_𝔄`$ there is an other seminorm $`\stackrel{~}{p}\mathrm{\Pi }_𝔄`$ such that for all $`X_1,X_2𝔄`$ the inequality $`p(X_1X_2)\stackrel{~}{p}(X_1)\stackrel{~}{p}(X_2)`$ holds. An inductive locally convex vector space $`𝔄`$ is an indexed family of locally convex vector spaces $`\{𝔄_\lambda \}_{\lambda \mathrm{\Lambda }}`$ such that the following holds: $`\mathrm{\Lambda }`$ is an upward directed partially ordered set, i. e. for all $`\lambda ,\mu \mathrm{\Lambda }`$ there is an element $`\nu \lambda ,\mu `$. For all $`\mu \lambda `$ there exist continuous inclusions $`T_\mu ^\lambda :𝔄_\lambda 𝔄_\mu `$; and for $`\nu \mu \lambda `$ one has $`T_\nu ^\mu T_\mu ^\lambda =T_\nu ^\lambda `$. Now, $`𝔄`$ is an inductive locally convex algebra if for each $`\lambda ,\mu \mathrm{\Lambda }`$ there is an element $`\mathrm{prod}(\lambda ,\mu )\mathrm{\Lambda }`$, and for $`\nu \mathrm{prod}(\lambda ,\mu )`$ bilinear products $`M_{\lambda ,\mu }^\nu :𝔄_\lambda \times 𝔄_\mu 𝔄_\nu `$ compatible with the inclusions and the usual algebraic prescriptions are given. An element of $`𝔄`$ is an element of $`_{\lambda \mathrm{\Lambda }}𝔄_\lambda `$ making identifications along the inclusion maps. Then $`𝔄`$ will be an algebra. endowed with an “inductive” topology coming from the filtration $`\{𝔄_\lambda \}_{\lambda \mathrm{\Lambda }}`$, such that the vector space structure respects the filtration but the algebra structure does not. If the spaces $`𝔄`$ and $`𝔅`$ have inductive topologies with filtrations $`\{𝔄_\lambda \}_{\lambda \mathrm{\Lambda }}`$ and $`\{𝔅_\mu \}_{\mu M}`$, then a map $`\varphi :𝔄𝔅`$ is continuous if for each $`\lambda \mathrm{\Lambda }`$ there is an element $`\mu M`$ such that there is a continuous map $`\varphi _\lambda :𝔄_\lambda 𝔅_\mu `$, which is set-theoretically a restriction of $`\varphi `$.
Suppose that $`\mathrm{\Theta }_1,\mathrm{\Theta }_2`$ are sets and $`𝔙`$ is a vector space. Then a $`𝔙`$-valued $`\mathrm{\Theta }_1`$ times $`\mathrm{\Theta }_2`$ matrix is just a formal sum $`s=_{a\mathrm{\Theta }_1,b\mathrm{\Theta }_2}s_{a,b}𝐞_{a,b}_{\mathrm{\Theta }_1,\mathrm{\Theta }_2}(𝔙)`$ with coefficients $`s_{a,b}`$ from $`𝔙`$. We write $`_\mathrm{\Theta }(𝔙)`$ instead of $`_{\mathrm{\Theta },\mathrm{\Theta }}(𝔙)`$ and use similar notation for other spaces as well. For column and row matrices, we use the notation $`𝐞_a=𝐞_{a,}`$ and $`𝐞_b^{}=𝐞_{,b}`$ respectively, and we make the formal identification $`𝐞_{a,b}=𝐞_a𝐞_b^{}`$. For column spaces, we use the notation $`𝒮(\mathrm{\Theta };𝔙)=_{\mathrm{\Theta },\{\}}(𝔙)`$. We use the notation $`1_\mathrm{\Theta }=_{\theta \mathrm{\Theta }}𝐞_{\theta ,\theta }`$, and in general circumstances we consider the identity matrix $`1_\mathrm{\Theta }`$ as the adjoint unit in any non-unital $`\mathrm{\Theta }`$ times $`\mathrm{\Theta }`$ matrix algebra. If $`s_i𝔙`$, $`i`$ are given then
(1)
$$\mathrm{Diag}(\mathrm{}s_2,s_1|s_0,s_1,s_2,\mathrm{})=\underset{i}{}s_i𝐞_{i,i}_{}(𝔙)$$
is the corresponding diagonal matrix. $`\mathrm{Diag}(s_0,s_1,s_2,\mathrm{})_{}(𝔙)`$, similarly. For $`a𝔄`$, we define the matrices $`𝖤_{}(a)=a𝐞_{00}𝒦_{}(𝔄)`$ and $`𝖤_{}(a)=a𝐞_{00}𝒦_{}(𝔄)`$. Then, as usual, for $`\stackrel{~}{a}=1+a1+𝔄`$, we extend these maps as $`𝖤_{}(\stackrel{~}{a})=1_{}+𝖤_{}(a)`$ and $`𝖤_{}(\stackrel{~}{a})=1_{}+𝖤_{}(a)`$; i. e., for $`\stackrel{~}{a}1+𝔄`$, it yields $`𝖤_{}(\stackrel{~}{a})=\mathrm{Diag}(\stackrel{~}{a},1,1,\mathrm{})`$, and $`𝖤_{}(\stackrel{~}{a})=\mathrm{Diag}(\mathrm{},1,1|\stackrel{~}{a},1,1,\mathrm{})`$.
On the set $``$ of natural numbers, there is the natural space $`𝒮^{\mathrm{}}(;)^{}`$, i. e. the space of multiplicatively invertible polynomially growing functions. A countable set $`\mathrm{\Theta }`$ is called a set of polynomial growth if it is endowed with a set of functions $`𝒮^{\mathrm{}}(\mathrm{\Theta };)^{}`$ from $`\mathrm{\Theta }`$ to $``$ such that there is a bijection $`\omega :\mathrm{\Theta }`$ such that $`\omega ^{}𝒮^{\mathrm{}}(;)^{}=𝒮^{\mathrm{}}(\mathrm{\Theta };)^{}`$. It is notable that $`\times `$ and $`\dot{}`$ are sets of polynomial growth naturally; and that way we can define the direct product $`\mathrm{\Theta }_1\times \mathrm{\Theta }_2`$ and direct sums $`\mathrm{\Theta }_1\dot{}\mathrm{\Theta }_2`$ of sets of polynomial growth $`\mathrm{\Theta }_1`$ and $`\mathrm{\Theta }_2`$. In what follows, the sets of polynomial growth we use will be like $`,,`$ or $`\{1,\mathrm{},n\}\times `$, where the description of the relevant function spaces is evident, so it will not be detailed. The main point is that a set $`\mathrm{\Theta }`$ of polynomial growth is just like $``$ for practical purposes. If $`\mathrm{\Theta }_1,\mathrm{\Theta }_2`$ are sets of polynomial growth, and $`𝔙`$ is a locally convex vector space, then we can define some matrix spaces as follows:
(a) With functions $`F:\mathrm{\Pi }_𝔙𝒮^{\mathrm{}}(\mathrm{\Theta }_1;)^{}\times 𝒮^{\mathrm{}}(\mathrm{\Theta }_2;)^{}`$, the filtering spaces
(2)
$$\begin{array}{c}_{\mathrm{\Theta }_1,\mathrm{\Theta }_2}^\mathrm{},\mathrm{}(𝔙)_F=\{s_{\mathrm{\Theta }_1,\mathrm{\Theta }_2}(𝔙):p\mathrm{\Pi }_𝔙\text{ }\hfill \\ \hfill |s|_{\frac{1}{F_1(p)},p,\frac{1}{F_2(p)}}=\underset{(a,b)\mathrm{\Theta }_1\times \mathrm{\Theta }_2}{}|\frac{1}{F_1(p)(a)}\left|p(s_{a,b})\right|\frac{1}{F_2(p)(b)}|<+\mathrm{}\}\end{array}$$
form the inductive locally convex space $`_{\mathrm{\Theta }_1,\mathrm{\Theta }_2}^\mathrm{},\mathrm{}(𝔙)`$.
(b) With functions $`F:\mathrm{\Pi }_𝔙\times 𝒮^{\mathrm{}}(\mathrm{\Theta }_2;)^{}𝒮^{\mathrm{}}(\mathrm{\Theta }_1;)^{}`$, the filtering spaces
(3)
$$\begin{array}{c}_{\mathrm{\Theta }_1,\mathrm{\Theta }_2}^\mathrm{},\mathrm{}(𝔙)_F=\{s_{\mathrm{\Theta }_1,\mathrm{\Theta }_2}(𝔙):p\mathrm{\Pi }_𝔙g𝒮^{\mathrm{}}(\mathrm{\Theta }_2;)^{}\text{ }\hfill \\ \hfill |s|_{\frac{1}{F(p,g)},p,g}=\underset{(a,b)\mathrm{\Theta }_1\times \mathrm{\Theta }_2}{}|\frac{1}{F(p,g)(a)}|p(s_{a,b})|g(b)|<+\mathrm{}\}\end{array}$$
form the inductive locally convex space $`_{\mathrm{\Theta }_1,\mathrm{\Theta }_2}^\mathrm{},\mathrm{}(𝔙)`$. We can define the space $`_{\mathrm{\Theta }_1,\mathrm{\Theta }_2}^\mathrm{},\mathrm{}(𝔙)`$ similarly.
(c) We define
(4)
$$\begin{array}{c}^\mathrm{},\mathrm{}(\mathrm{\Theta }_1,\mathrm{\Theta }_2;𝔙)=\{s(\mathrm{\Theta }_1,\mathrm{\Theta }_2;𝔙):p\mathrm{\Pi }_𝔙f𝒮^{\mathrm{}}(\mathrm{\Theta }_1;)^{}\hfill \\ \hfill g𝒮^{\mathrm{}}(\mathrm{\Theta }_2;)^{}\text{ }|s|_{f,p,g}=\underset{(a,b)\mathrm{\Theta }_1\times \mathrm{\Theta }_2}{}|f(a)|p(s_{a,b})|g(b)|<+\mathrm{}\}.\end{array}$$
(d) It is natural to define $`\mathrm{\Psi }_{\mathrm{\Theta }_1,\mathrm{\Theta }_2}(𝔙)=_{\mathrm{\Theta }_1,\mathrm{\Theta }_2}^\mathrm{},\mathrm{}(𝔙)_{\mathrm{\Theta }_1,\mathrm{\Theta }_2}^\mathrm{},\mathrm{}(𝔙),`$ the space of matrices of “pseudodifferential size”.
If $`𝔄\times 𝔅`$ is a continuous bilinear pairing between locally convex spaces, then we have induced continuous pairings $`_{\mathrm{\Theta }_1,\mathrm{\Theta }_2}^{𝖷,\mathrm{}}(𝔄)\times _{\mathrm{\Theta }_2,\mathrm{\Theta }_3}^{\mathrm{},𝖸}(𝔅)_{\mathrm{\Theta }_1,\mathrm{\Theta }_3}^{𝖷,𝖸}()`$, $`\mathrm{\Psi }_{\mathrm{\Theta }_1,\mathrm{\Theta }_2}(𝔄)\times \mathrm{\Psi }_{\mathrm{\Theta }_2,\mathrm{\Theta }_3}(𝔅)\mathrm{\Psi }_{\mathrm{\Theta }_1,\mathrm{\Theta }_3}(),`$ etc. So come the algebra and module structures associated to matrices. Instead of $`_{\mathrm{\Theta }_1,\mathrm{\Theta }_2}^\mathrm{},\mathrm{}(𝔄)`$ we will use the shorter notations $`𝒦_{\mathrm{\Theta }_1,\mathrm{\Theta }_2}(𝔄)`$. Instead of $`_{\mathrm{\Theta },\{\}}^{\pm \mathrm{},𝖷}(𝔄)`$ it is reasonable to use $`𝒮^\pm \mathrm{}(\mathrm{\Theta };𝔄)`$, which is a consistent extension of our earlier notation. There are natural isomorphisms like $`𝒦_{\mathrm{\Theta }_1\times \mathrm{\Theta }_1^{},\mathrm{\Theta }_2\times \mathrm{\Theta }_2^{}}(𝔄)𝒦_{\mathrm{\Theta }_1,\mathrm{\Theta }_2}(𝒦_{\mathrm{\Theta }_1^{},\mathrm{\Theta }_2^{}}(𝔄))`$, etc. One often uses is relabeling of matrices, which is as follows: Suppose that $`\omega :\mathrm{\Omega }\mathrm{\Omega }^{}`$ is a map between sets of polynomial growth, such that $`\omega ^{}𝒮^{\mathrm{}}(\mathrm{\Omega }^{},)^{}=𝒮^{\mathrm{}}(\mathrm{\Omega },)^{}`$. This includes the case when $`\omega `$ is an isomorphism of sets of polynomial growth, and also the natural inclusions $`\iota :\mathrm{\Omega }\mathrm{\Omega }^{}\dot{}\mathrm{\Omega }^{\prime \prime }`$, where $`\mathrm{\Omega }^{\prime \prime }`$ is finite or an other set of polynomial growth. Let us now consider the matrix $`R_\omega =_{\alpha \mathrm{\Omega }}𝐞_{\alpha ,\omega (\alpha )}\mathrm{\Psi }_{\mathrm{\Omega },\mathrm{\Omega }^{}}()`$. Then, for a matrix $`A_\mathrm{\Theta }^{𝖷,𝖸}(𝔄)`$ or $`\mathrm{\Psi }_\mathrm{\Theta }(𝔄)`$, we can take the matrix $`r_\omega (A)=R_\omega ^{}AR_\omega `$ which is a matrix of the same kind as $`A`$ but $`\mathrm{\Omega }`$ is replaced by $`\mathrm{\Omega }^{}`$. This relabeling $`r_\omega `$ is a continuous, smooth operation, which is an isomorphism if $`\omega `$ is an isomorphism.
The advantage of the spaces $`_{\mathrm{\Theta }_1,\mathrm{\Theta }_2}^{𝖷,𝖸}(𝔙)`$ is that they are sufficiently large for the purposes of arithmetic calculations. In what follows, only the algebras $`𝒦`$ will be used explicitly. On the other hand, all computations, except in Section 7 will be governed by the principle every matrix expression will be understood as an element of $`\mathrm{\Psi }_{\mathrm{\Omega }_1,\mathrm{\Omega }_2}(𝔄)`$, where $`\mathrm{\Omega }_i`$ are sets of polynomial growth, and $`𝔄`$ is a locally convex algebra; but we always hope that our expressions will yield results which turn out to be continuous in stronger topologies.
## 3. The environment of cyclic and Toeplitz algebras
### Cyclic and Toeplitz algebras
In what follows, let $`\overline{}=`$, so $`=\overline{}\dot{}`$. We make a canonical correspondence between $``$ and $`\overline{}`$ by relabeling every $`n`$ to $`1n`$. We can consider every $`\times `$ matrix $`U`$ as a $`2\times 2`$ matrix of $`\times `$ matrices:
$$U=\left[\begin{array}{cc}U|_{\overline{}\times \overline{}}& U|_{\overline{}\times }\\ & \\ U|_{\times \overline{}}& U|_\times \end{array}\right]\left[\begin{array}{cc}U^{}& U^+\\ U^+& U^{++}\end{array}\right],$$
such that the matrix entries on the right side are $`\times `$ matrices obtained by the correspondence explained above.
An element $`a=_ia_i𝐞_i𝒮^{\mathrm{}}(;𝔄)`$ can and will, in general, be identified with the Laurent series $`_ia_i𝗓^i𝔄[𝗓^1,𝗓]`$ with rapidly decreasing coefficients. We call this algebra the algebra of cyclic loops, in contrast to the algebra of proper loops $`𝒞^{\mathrm{}}(\mathrm{S}^1;𝔄)`$. Elements $`a=_ia_i𝗓^i𝔄[𝗓^1,𝗓]`$ can be represented by $`\times `$ matrices
(5)
$$𝖴(a)=\underset{n,m}{}a_{nm}𝐞_{n,m}=\left[\begin{array}{cccccc}\mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \mathrm{}& a_0& a_1& a_2& a_3& \mathrm{}\\ \mathrm{}& a_1& a_0& a_1& a_2& \mathrm{}\\ & & & & & \\ \mathrm{}& a_2& a_1& a_0& a_1& \mathrm{}\\ \mathrm{}& a_3& a_2& a_1& a_0& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right].$$
If $`𝒲_{}(𝔄)`$ is the image set of $`𝔄[𝗓^1,𝗓]`$ under $`𝖴`$, then it is a subset of $`\mathrm{\Psi }(;𝔙)`$ algebraically, but we put the topology of $`𝒮^{\mathrm{}}(;𝔙)`$ to it. If $`𝔄`$ is a locally convex algebra, then $`𝖴:𝔄[𝗓^1,𝗓]𝒲_{}(𝔄)`$ is an isomorphism of algebras. When it comes to the $`2\times 2`$ decomposition as explained above, in order to simplify the notation, we will just write $`𝖶(a)`$ instead of $`𝖴(a)^{++}`$, and $`𝖸(a)`$ instead of $`𝖴(a)^+`$. If $`a=a(𝗓)`$ then, with some abuse of notation, we also write $`a^{}=a(𝗓^1)`$. Then
$$𝖴(a)=\left[\begin{array}{cc}𝖶(a^{})& 𝖸(a^{})\\ 𝖸(a)& 𝖶(a)\end{array}\right].$$
So $`𝖶(a)`$ is the infinite Toeplitz matrix associated to $`a`$, and $`𝖸(a)`$ is the infinite Hankel matrix associated to (”the positive part” of) $`a`$.
As far as the linear structure is concerned, we could have just used the matrices $`𝖶(a)`$ to represent the elements $`a`$. The difference is that in terms of matrix multiplication $`𝖶(a)𝖶(b)=𝖶(ab)𝖸(a)𝖸(b^{}),`$ so there is an “anomalous” term $`𝖸(a)𝖸(b^{})𝒦_{}(𝔄)`$. One can see that algebraically $`𝒲_{}(𝔄)𝒦_{}(𝔄)=0`$. Hence it is reasonable to define the Toeplitz algebra
$$𝒯_{}(𝔄)=𝒲_{}(𝔄)+𝒦_{}(𝔄),$$
which is topologically just $`𝒲_{}(𝔄)𝒦_{}(𝔄)`$ but with the algebraic product rule $`(𝖶(a)+p)((𝖶(b)+q))=𝖶(ab)+(𝖸(a)𝖸(b^{})+𝖶(a)q+p𝖶(b)+pq),`$ induced from the matrix structure. Algebraically, $`𝒯_{}(𝔄)`$ is just a subset of $`\mathrm{\Psi }(;𝔄)`$ but a locally convex algebra. So, one can see that there is a short exact sequence of algebras $`0𝒦_{}(𝔄)\stackrel{𝜄}{}𝒯_{}(𝔄)\stackrel{𝜎}{}𝒲_{}(𝔄)0`$. The map $`\iota `$ is the inclusion of the ideal of rapidly decreasing matrices into the Toeplitz algebra, while $`\sigma `$ is the symbol map. In what follows, we rather consider the value of the symbol map as an element of $`𝔄[𝗓^1,𝗓]`$, so we have the symbol homomorphism
$$\sigma :𝒯_{}(𝔄)𝔄[𝗓^1,𝗓].$$
We can naturally extend this symbol map to unit groups as we have seen.
For technical reasons, we define the algebra
$$𝒯_{}(𝔄)=\left[\begin{array}{cc}𝒯_{}(𝔄)& 𝒦_{}(𝔄)\\ 𝒦_{}(𝔄)& 𝒯_{}(𝔄)\end{array}\right],$$
which is also naturally a locally convex algebra. Then $`𝒲_{}(𝔄)𝒲_{}(𝔄)+𝒦_{}(𝔄)𝒯_{}(𝔄)`$. For the sake of notational convenience, we define the block matrix
$$\widehat{𝖴}(a)=\left[\begin{array}{ccc}𝖶(a^{})& & 𝖸(a^{})\\ & & \\ & 0& \\ 𝖸(a)& & 𝖶(a)\end{array}\right]𝒯_{}(𝔄).$$
We remark that for $`\stackrel{~}{a}𝚄(𝔄[𝗓^1,𝗓])`$ an “$`1`$” appears in the place of “$`0`$”.
Elements of $`𝒯_{}(𝔄)`$ have two symbols; one belonging to the lower right quadrant, and one belonging to the upper left quadrant. It is a small but important observation regarding $`𝖴(a)𝒯_{}(𝔄)`$ that the Toeplitz element in the lower right quadrant has symbol $`a=a(𝗓)`$, but the Toeplitz element in the upper left quadrant has symbol $`a^{}=a(𝗓^1)`$.
One can also see that there are natural isomorphisms like $`𝒯_{}(𝒦_\mathrm{\Omega }(𝔄))𝒦_\mathrm{\Omega }(𝒯_{}(𝔄)),`$ etc. In fact, all of our matrix space constructions considered as functors are naturally “commutative”.
Let $`𝔄[𝗓^1,𝗓]^{\mathrm{po}}`$ be the set of pointed loops, i. e., where $`a(1)=0`$. Then the elements $`\stackrel{~}{a}𝚄(𝔄[𝗓^1,𝗓]^{\mathrm{po}}`$) are those for which $`\stackrel{~}{a}(1)=1`$. These pointed spaces are closed subspaces of the unpointed spaces. We can define the pointed Toeplitz algebra $`𝒯_{}(𝔄)^{\mathrm{po}}`$ similarly, the symbols are pointed there.
### The Bott involution map
In what follows, we use the abbreviation $`\mathrm{\Lambda }(a,b)=\frac{1}{2}(1+a+bab)`$. Let $`𝖰=\left[\begin{array}{cc}1_{}& \\ & 1_{}\end{array}\right]`$. Then $`\mathrm{\Lambda }(𝗓,𝖰)=\left[\begin{array}{cc}𝗓1_{}& \\ & 1_{}\end{array}\right].`$ We also use the delta-function $`\delta _{n,m}`$, which is 1 if $`n=m`$, and it is $`0`$ otherwise.
If $`a𝚄(𝔄[𝗓^1,𝗓])`$, then we define the “Bott” involution
$$𝖡(a)=𝖴(a)\mathrm{𝖰𝖴}(a)^1𝖰+𝒦_{}(𝔄)$$
(cf. the symbols).
### “Shifting rotations”
Our natural deformation parameter variable, in general, will be $`\theta [0,\pi /2]`$, or, more generally, $`\theta \mathrm{S}^1=/2\pi `$. In order to save space, we often use $`t=\mathrm{sin}\theta `$ and $`s=\mathrm{cos}\theta `$ instead. It is useful to keep in mind that $`s^2=1t^2`$. For $`\theta \mathrm{S}^1`$, we define the matrices
$$𝖢(\theta )=\left[\begin{array}{ccccc}s& ts& t^2s& t^3s& \mathrm{}\\ t& s^2& ts^2& t^2s^2& \mathrm{}\\ & t& s^2& ts^2& \mathrm{}\\ & & t& s^2& \mathrm{}\\ & & & t& \mathrm{}\\ & & & & \mathrm{}\end{array}\right]𝒯_{}().$$
###### Lemma 3.1.
Let $`𝖢(\theta )^{}`$ denote the transpose of $`𝖢(\theta )`$. Then
(a)
$$𝖢(\theta )^{}𝖢(\theta )=1_{}.$$
(b)
$$𝖢(\theta )𝖢(\theta )^{}=\delta _{t,1}𝐞_{0,0}\delta _{t,1}𝐞_{0,0}+1_{}.$$
(c)
$$𝖢(\theta )𝐞_{n,m}𝖢(\theta )^{}=\left[\begin{array}{cccccc}t^{n+m}s^2& t^{n+m1}s^3& \mathrm{}& t^ns^3& t^{n+1}s& \\ t^{n+m1}s^3& t^{n+m2}s^4& \mathrm{}& t^{n1}s^4& t^ns^2& \\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ t^ms^3& t^{m1}s^4& \mathrm{}& s^4& ts^2& \\ t^{m+1}s& t^ms^2& \mathrm{}& ts^2& t^2\\ & & & & & 0\\ & & & & & & \mathrm{}\end{array}\right].$$
(d) For $`n>0`$,
$$𝖢(\theta )𝖶(𝗓^n)𝖢(\theta )^{}=\delta _{t,1}𝐞_{0,0}(1)^n\delta _{t,1}𝐞_{0,0}+\left[\begin{array}{cc}t^n& \\ t^{n1}s& \\ \mathrm{}\\ ts& \\ s& \\ & 1\\ & & 1\\ & & & \mathrm{}\end{array}\right],$$
and $`𝖢(\theta )𝖶(𝗓^n)𝖢(\theta )^{}`$ is the transpose of the matrix above.
###### Proof.
Direct computation. ∎
###### Remark 3.2.
The presence of the terms $`\delta _{t,1},\delta _{t,1}`$ might be surprising at first sight. It reflects the phenomenon that in a topological algebra one cannot simultaneously topologize the families $`𝖢(\theta )`$ and $`𝖢(\theta )^{}`$ correctly. In fact, the $`𝖢(\theta )`$’s are isometries for $`1<t<1`$, but they are just partial isometries for $`t=\pm 1`$.
## 4. Stabilizing homotopies
###### Proposition 4.1.
The continuous map
$$T_𝒦:𝒦_{}(𝔄)\times \mathrm{S}^1𝒦_{}(𝔄)$$
given by
$$A,\theta T_𝒦(A,\theta )=𝖢(\theta )A𝖢(\theta )^{}$$
is smooth in $`\theta `$. It yields a family of endomorphisms of $`𝒦_{}(𝔄)`$ when $`\theta `$ is fixed. These are isomorphisms for $`1<t<1`$, and closed injective endomorphisms for $`t=\pm 1`$. In particular, for $`\theta =0`$ $`(t=0)`$,
$$T_𝒦(A,0)=A=\left[\begin{array}{ccc}a_{11}& a_{12}& \mathrm{}\\ a_{21}& a_{22}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right];$$
but for $`\theta =\pm \pi /2`$ ($`t=\pm 1`$),
$$T_𝒦(A,\pm \pi /2)=\left[\begin{array}{c}0\\ & a_{11}& a_{12}& \mathrm{}\\ & a_{21}& a_{22}& \mathrm{}\\ & \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right].$$
###### Proof.
Well-definedness and smoothness follows from Lemma 3.1.c. Lemma 3.1.a implies that we have a family of endomorphisms. Furthermore, it also shows that $`A=𝖢(\theta )^{}T_𝒦(A,\theta )𝖢(\theta )`$; from which the statement about the nature of the endomorphisms follows easily. ∎
Hence, taking $`\theta [0,\pi /2]`$, we see that the deformation $`T_𝒦`$ does indeed realize a stabilizing homotopy, even if only with one “extra dimension”. Nevertheless, after this, stabilization becomes a matter of standard tricks:
###### Corollary 4.2 ($``$ Statement 1.1).
Let $`\mathrm{\Omega }_1`$ and $`\mathrm{\Omega }_2`$ be sets of polynomial growth; $`\mathrm{\Omega }=\mathrm{\Omega }_1\dot{}\mathrm{\Omega }_2`$, and let $`\omega :\mathrm{\Omega }\mathrm{\Omega }_1\mathrm{\Omega }`$ be the composition of an isomorphism $`\mathrm{\Omega }\mathrm{\Omega }_1`$ and the natural inclusion $`\mathrm{\Omega }_1\mathrm{\Omega }_1\dot{}\mathrm{\Omega }_2`$. Then we claim:
There is a smooth map $`\widehat{T}_𝒦:𝒦_\mathrm{\Omega }(𝔖)\times [0,\pi /2]𝒦_\mathrm{\Omega }(𝔖)`$ such that it yields a family of endomorphisms of $`𝒦_\mathrm{\Omega }(𝔖)`$, which are isomorphisms for $`\theta [0,\pi /2)`$, and a closed injective endomorphism for $`\theta =\pi /2`$, such that $`\widehat{T}_𝒦(A,0)=\mathrm{id}_{𝒦_\mathrm{\Omega }(𝔖)}`$ and $`\widehat{T}_𝒦(A,\pi /2)=r_\omega `$. The map $`\widehat{T}_𝒦`$ extends to unit groups naturally.
###### Proof.
Take $`𝔄=𝒦_{}(𝔖)`$ in the previous statement. It yields our statement with $`\mathrm{\Omega }_1=(\{0\})\times `$, $`\mathrm{\Omega }_2=\{0\}\times `$, $`\omega ((n,m))=(n+1,m)`$. Now, using an appropriate relabeling $`r_\eta `$ of $`\mathrm{\Omega }`$ we obtain the general statement. ∎
###### Remark 4.3.
Another way to achieve stabilization by many dimensions is to “quantize” $`𝖢(\theta )`$, see .
Due to the multiplicative structure, the concatenation of group valued homotopies is particularly simple: If $`f,g:Y\times [0,1]G`$ yield homotopies $`f_0f_1`$, $`g_0g_1`$ where $`f_1=g_0`$, then $`h(y,t)=f(y,t)f(y,1)^1g(y,t)`$ yields a homotopy between $`f_0`$ and $`g_1`$. Then polynomial/smooth homotopies yield polynomial/smooth homotopies, and the operation is associative; in contrast to concatenation by reparametrization. Using this observation and the stabilizing homotopies above, one can easily prove
###### Corollary 4.4.
Let $`\mathrm{\Omega }_1,\mathrm{\Omega }_2`$ be sets of polynomial growth, and let $`\iota _1:\mathrm{\Omega }_1\mathrm{\Omega }_1\dot{}\mathrm{\Omega }_2`$ be the natural inclusion. Assume that $`H:X\times [0,1]𝚄(𝒦_{\mathrm{\Omega }_1\dot{}\mathrm{\Omega }_2}(𝔖))`$ is a smooth homotopy with maps $`f_0,f_1:X𝚄(𝒦_{\mathrm{\Omega }_1}(𝔖))`$ such that $`H_0=r_{\iota _1}(f_0)`$ and $`H_1=r_{\iota _1}(f_1)`$. Then we claim that there is a smooth homotopy $`f:X\times [0,1]𝚄(𝒦_{\mathrm{\Omega }_1}(𝔖))`$ between $`f_0`$ and $`f_1`$. This $`f`$ can be chosen so that there is a smooth homotopy between $`H`$ and $`r_{\iota _1}(f)`$ relative to endpoints. In other words: “In stable algebras stable homotopies can be reduced to ordinary homotopies.” ∎
## 5. Linearization of cyclic loops
###### 5.1.
Let $`𝗏`$ be a cyclic formal variable, and take $`𝖵=_n𝗏^n𝐞_{n,n}`$. Furthermore, take $`𝖦(\theta ,𝗏)=\left[\begin{array}{cc}1_{}& \\ & 𝖵^1𝖢(\theta )𝖵\end{array}\right]`$ and $`𝖦^{}(\theta ,𝗏)=\left[\begin{array}{cc}1_{}& \\ & 𝖵^1𝖢(\theta )^{}𝖵\end{array}\right].`$ For $`a𝔄[𝗓^1,𝗓]`$, we define
$$U(a,\theta ,𝗏)=\delta _{t,1}a(𝗏)𝐞_{00}+\delta _{t,1}a(𝗏)𝐞_{00}+𝖦(\theta ,𝗏)𝖴(a)𝖦^{}(\theta ,𝗏).$$
###### 5.2.
For $`n>0`$, this definition yields
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$$\begin{array}{c}U(𝗓^n,\theta ,𝗏)=\hfill \\ \hfill =\left[\begin{array}{cccccccccc}\mathrm{}& & & & & & & & & \\ & 1& & & & & & & & \\ & & & & & & & & & \\ & & s& ts𝗏& t^2s𝗏^2& \mathrm{}& t^{n1}s𝗏^{n1}& t^n𝗏^n& & \\ & & t𝗏^1& s^2& ts^2𝗏& \mathrm{}& t^{n2}s^2𝗏^{n2}& t^{n1}s𝗏^{n1}& & \\ & & & \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& & \\ & & & & \mathrm{}& \mathrm{}& ts^2𝗏& t^2s𝗏^2& & \\ & & & & & \mathrm{}& s^2& ts𝗏& & \\ & & & & & & t𝗏^1& s& & \\ & & & & & & & & 1& \\ & & & & & & & & & \mathrm{}\end{array}\right]\end{array}$$
$`U(1,\theta ,𝗏)=1_{},`$ and $`U(𝗓^n,\theta ,𝗏)=U(𝗓^n,\theta ,𝗏^1)^{}`$; i. e. $`U(𝗓^n,\theta ,𝗏)`$ is just a rather nice perturbation / deformation of $`𝖴(𝗓^n)`$.
###### Lemma 5.3.
The continuous map
$$U:𝔄[𝗓^1,𝗓]\times \mathrm{S}^1𝒲_{}(𝔄)+𝒦_{}(𝔄)[𝗏^1,𝗏]$$
defined by
$$a,\theta U(a,\theta ,𝗏)$$
is smooth in $`\theta `$. It yields a family of homomorphisms with fixed $`\theta `$. The symbols remain constant. For $`\theta =0`$,
$$U(a,0,𝗏)=\left[\begin{array}{cc}𝖶(a^{})& 𝖸(a^{})\\ 𝖸(a)& 𝖶(a)\end{array}\right]=𝖴(a);$$
while for $`\theta =\pi /2`$,
$$U(a,\pi /2,𝗏)=\left[\begin{array}{ccc}𝖶(a^{})& & \mathrm{𝗏𝖸}(a^{})\\ & & \\ & a(𝗏)& \\ 𝗏^1𝖸(a)& & 𝖶(a)\end{array}\right].$$
###### Proof.
This is immediate from 5.2 by taking linear combinations. ∎
Considering $`a𝚄(𝔄)`$, and the natural extension to the unit group, $`U(a,\pi /2,𝗏)=𝖤_{}(a(𝗏))\mathrm{\Lambda }(𝗏,𝖰)\widehat{𝖴}(a)\mathrm{\Lambda }(𝗏,𝖰)^1`$ can be written.
###### Proposition 5.4 ($``$ Statement 1.2).
The continuous map
$$K:𝚄(𝔄[𝗓^1,𝗓])\times \mathrm{S}^1𝚄(𝒦_{}(𝔄)[𝗏^1,𝗏]^{\mathrm{po}})$$
defined by
$$a,\theta K(a,\theta ,𝗏)=U(a,\theta ,𝗏)\mathrm{\Lambda }(𝗏,𝖰)U(a,\theta ,1)^1\mathrm{\Lambda }(𝗏,𝖰)^1$$
is smooth in the variable $`\theta `$. Here
$$K(a,0,𝗏)=\mathrm{\Lambda }(𝗏,𝖡(a))\mathrm{\Lambda }(𝗏,𝖰)^1,K(a,\pi /2,𝗏)=𝖤_{}(a(𝗏)a(1)^1).$$
###### Proof.
The statement follows immediately from the previous lemma.∎
###### Remark 5.5.
When it comes to the linearization of not pointed loops but the “cocycle” $`a(𝗓)a(𝗐)^1`$, then one can use the linearizing “cocycle” $`K^\mathrm{c}(a,\theta ,𝗓,𝗐)=U(a,\theta ,𝗓)\mathrm{\Lambda }(\mathrm{𝗓𝗐}^1,𝖰)`$ $`U(a,\theta ,𝗐)^1`$. It yields $`K^\mathrm{c}(a,0,𝗓,𝗐)=\mathrm{\Lambda }(\mathrm{𝗓𝗐}^1,𝖡(a))`$ and $`K^\mathrm{c}(a,\pi /2,𝗓,𝗐)=𝖤_{}(a(𝗓)a(𝗐)^1)`$. Then $`K(a,\theta ,𝗏)=K^\mathrm{c}(a,\theta ,𝗓,1)\mathrm{\Lambda }(𝗓,𝖰)^1`$.
It is notable that loops which are already linear will remain constant but stabilized: If $`a(𝗓)=\mathrm{\Lambda }(𝗓,\stackrel{~}{Q})`$ then $`K^\mathrm{c}(a,\theta ,𝗓,𝗐)=\mathrm{Diag}(\mathrm{},\mathrm{𝗓𝗐}^1|\mathrm{\Lambda }(\mathrm{𝗓𝗐}^1,\stackrel{~}{Q}),1\mathrm{})`$, independently from $`\theta `$. Similarly, rapidly decreasing perturbations of a linear loop will linearize through rapidly decreasing perturbations of that linear loop.
###### Remark 5.6.
For a locally convex algebra $`𝔄`$ we can define
$$K_0(𝔄)=\pi _0^{\mathrm{smooth}}(𝔫𝔳𝔬𝔩(𝖰+𝒦_{}(𝔄))),$$
the smooth path components of the involutions, which are perturbations of $`𝖰`$. Similarly, one can define
$$K_1(𝔄)=\pi _0^{\mathrm{smooth}}(𝚄(𝒦_{}(𝔄))).$$
Now $`𝖡`$, by this linearization argument, induces an isomorphism
$$𝖡_{}:K_1(𝔄[𝗓^1,𝗓]^{\mathrm{po}})K_0(𝔄).$$
This is the “hard part” of Bott periodicity in the complex case, when geometric loops can be represented by cyclic loops.
## 6. The contractibility of the pointed stable Toeplitz unit group
When we extend the stabilization procedure of Proposition 4.1 to Toeplitz algebras, the symbol suddenly appears in the result:
###### Proposition 6.1.
The continuous map
$$T:𝒯_{}(𝔄)\times \mathrm{S}^1𝒲_{}(𝔄)+𝒦_{}(𝔄)[𝗏^1,𝗏]𝒯_{}(𝔄)[𝗏^1,𝗏]$$
defined by
$$A,\theta T(A,\theta ,𝗏)=\delta _{t,1}a(𝗏)𝐞_{0,0}+\delta _{t,1}a(𝗏)𝐞_{0,0}+𝖵^1𝖢(\theta )𝖵A𝖵^1𝖢(\theta )^{}𝖵,$$
where $`a=\sigma (A)`$, is smooth in the variable $`\theta `$. It yields a family of homomorphisms of $`𝒯_{}(𝔄)`$ to $`𝒯_{}(𝔄)[𝗏^1,𝗏]`$. The map leaves the symbol invariant. For $`\theta =0`$,
$$T(A,0,𝗏)=A=\left[\begin{array}{ccc}a_{11}& a_{12}& \mathrm{}\\ a_{21}& a_{22}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right];$$
but for $`\theta =\pm \pi /2`$,
$$T(A,\pm \pi /2,𝗏)=\left[\begin{array}{c}a(\pm 𝗏)\\ & a_{11}& a_{12}& \mathrm{}\\ & a_{21}& a_{22}& \mathrm{}\\ & \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right].$$
###### Proof.
It follows from direct inspection of the matrices in question. ∎
As a corollary we obtain
###### Proposition 6.2.
The continuous map $`Z:𝚄(𝒯_{}(𝔄))\times \mathrm{S}^1𝚄(𝒦_{}(𝔄)[𝗏^1,𝗏]^{\mathrm{po}})`$ defined by $`A,\theta Z(A,\theta ,𝗏)=T(A,\theta ,𝗏)T(A,\theta ,1)^1`$ is smooth in $`\theta `$. For $`\theta =0`$ it yields $`Z(A,0,𝗏)=1_{},`$ but for $`\theta =\pm \pi /2`$ it yields $`Z(A,\pm \pi /2,𝗏)=𝖤_{}(a(\pm 𝗏)a(1)^1)`$.
Consequently, the symbols $`a(𝗓)`$ of invertible Toeplitz algebra elements are stably homotopic to constant loops $`a(1)`$. If $`𝔄=𝒦_{}(𝔖)`$, then (according to Corollary 4.2) stable homotopy implies the existence of ordinary homotopies. ∎
###### 6.3.
Suppose that $`Q`$ is an involution, and $`k𝔄`$. We will use the shorthand notation $`k_Q^+=\frac{1}{2}(k+QkQ)`$, $`k_Q^{++}=\frac{1+Q}{2}k\frac{1+Q}{2}`$, $`k_Q^+=\frac{1+Q}{2}k\frac{1Q}{2}`$, $`k_Q^+=\frac{1Q}{2}k\frac{1+Q}{2}`$. Let us define
$$L(Q,k)=𝖶(\mathrm{\Lambda }(𝗓,Q))k𝖶(\mathrm{\Lambda }(𝗓^1,Q))=\left[\begin{array}{cc}k_Q^{++}& k_Q^+\\ k_Q^+& k_Q^+& k_Q^+\\ & k_Q^+& k_Q^+& \mathrm{}\\ & & \mathrm{}& \mathrm{}\end{array}\right].$$
This is a homomorphism in $`k`$, and we can extend it to $`\stackrel{~}{k}=1+k`$ by $`L(Q,\stackrel{~}{k})=1_{}+L(Q,k)`$. Notice that in this case, $`\stackrel{~}{k}L(Q,\stackrel{~}{k}^1)`$ has symbol $`\stackrel{~}{k}\mathrm{\Lambda }(𝗓,Q)\stackrel{~}{k}^1\mathrm{\Lambda }(𝗓,Q)^1`$.
###### 6.4.
Assume that $`Q=𝖰`$ and $`k𝒯_{}(𝔖)`$. Set
$$\stackrel{~}{L}(k)=\left[\begin{array}{ccccccc}\mathrm{}& \mathrm{}& & & & & \\ \mathrm{}& k_Q^+& k_Q^+& & & & \\ & k_Q^+& k_Q^+& k_Q^+𝐞_{00}& k_Q^+𝐞_{10}& k_Q^+𝐞_{20}& \mathrm{}\\ & & & & & & \\ & & 𝐞_{00}k_Q^+& 𝐞_{00}k_Q^{++}𝐞_{00}& 𝐞_{00}k_Q^{++}𝐞_{10}& 𝐞_{00}k_Q^{++}𝐞_{20}& \mathrm{}\\ & & 𝐞_{01}k_Q^+& 𝐞_{01}k_Q^{++}𝐞_{00}& 𝐞_{01}k_Q^{++}𝐞_{10}& 𝐞_{01}k_Q^{++}𝐞_{20}& \mathrm{}\\ & & 𝐞_{02}k_Q^+& 𝐞_{02}k_Q^{++}𝐞_{00}& 𝐞_{02}k_Q^{++}𝐞_{10}& 𝐞_{02}k_Q^{++}𝐞_{20}& \mathrm{}\\ & & \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right].$$
What happens here, compared to $`L(Q,k)`$, is the following: We inflated the first row and column to infinitely many rows and columns, and reordered the matrix. Again, this is a homomorphism in $`k`$, and we can extend it to $`\stackrel{~}{k}1_{}+𝒯_{}(𝔖)`$ by taking $`\stackrel{~}{L}(\stackrel{~}{k})=1_\times +\stackrel{~}{L}(k)`$. Assume now that $`\stackrel{~}{k}𝚄(𝒯_{}(𝔖))`$, and the symbol of its lower right quadrant is $`a(𝗓)`$. Consider
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$$\begin{array}{c}𝖴(\mathrm{\Lambda }(𝗓,𝖰))\stackrel{~}{L}(\stackrel{~}{k})\mathrm{\Lambda }(\stackrel{~}{k}^1,𝖰_{𝒯_{}(𝔖)})𝖴(\stackrel{~}{k}\mathrm{\Lambda }(𝗓,𝖰)^1\stackrel{~}{k}^1)=\hfill \\ \hfill =\left[\begin{array}{ccccc}\mathrm{}& \mathrm{}& & & \\ & \frac{1𝖰}{2}& \frac{1+𝖰}{2}& & \\ & & & & \\ & & \frac{1𝖰}{2}& \frac{1+𝖰}{2}& \\ & & & \mathrm{}& \mathrm{}\end{array}\right]\stackrel{~}{L}(\stackrel{~}{k})\left[\begin{array}{cccccc}\mathrm{}& \mathrm{}& & & & \\ & \frac{1+𝖰}{2}& \frac{1𝖰}{2}& & & \\ & & & & & \\ & & \stackrel{~}{k}\frac{1+𝖰}{2}& \stackrel{~}{k}\frac{1𝖰}{2}& & \\ & & & \mathrm{}& \mathrm{}& \end{array}\right]\stackrel{~}{k}^1.\end{array}$$
From the observation $`\stackrel{~}{L}(\stackrel{~}{k})\mathrm{\Lambda }(\stackrel{~}{k}^1,𝖰_{𝒯_{}(𝔖)})𝚄(𝒯_{}(𝒦_{}(𝔖)))`$, and a careful examination of the matrix product, we find that the resulting expression is of shape $`\left[\begin{array}{cc}1_{}& \\ & N(\stackrel{~}{k})\end{array}\right]𝚄(𝒯_{}(𝒦_{}(𝔖)))`$; where we introduced the notation $`N(\stackrel{~}{k})`$ for the lower right quadrant. Then the component $`N(\stackrel{~}{k})𝚄(𝒯_{}(𝒦_{}(𝔖)))`$ has symbol $`\mathrm{\Lambda }(𝗓,𝖰)𝖤_{}(a(𝗓))\stackrel{~}{k}\mathrm{\Lambda }(𝗓,𝖰)^1\stackrel{~}{k}^1=𝖤_{}(a(𝗓))\mathrm{\Lambda }(𝗓,𝖰)\stackrel{~}{k}\mathrm{\Lambda }(𝗓,𝖰)^1\stackrel{~}{k}^1`$. Let us set $`G(a)=N(𝖴(a))`$. This yields
###### Proposition 6.5.
The continuous map $`G:𝚄(𝔖[𝗓^1,𝗓])𝚄(𝒯_{}(𝒦_{}(𝔖)))`$ is such that the symbol of $`G(a)`$ is $`𝖤_{}(a(𝗓))\mathrm{\Lambda }(𝗓,𝖰)𝖴(a)\mathrm{\Lambda }(𝗓,𝖰)^1𝖴(a)^1.`$
Now, according to Proposition 6.2 and Corollary 4.4, the mere existence of the map above implies that the symbol $`𝖤_{}(a(𝗓))\mathrm{\Lambda }(𝗓,𝖰)𝖴(a)\mathrm{\Lambda }(𝗓,𝖰)^1𝖴(a)^1`$ is homotopic to $`1_{}`$ for $`a(𝗓)𝚄(𝔄[𝗓^1,𝗓]^{\mathrm{po}})`$. So, Proposition 6.5 can be considered as a reformulation of linearizability.
###### Proposition 6.6 ($``$ Statement 1.3).
The unit group $`𝚄(𝒯_{}(𝒦_{}(𝔖))^{\mathrm{po}})`$ is smoothly contractible.
###### Proof.
We prove the statement up to stabilization. Then stabilization can be removed according to Corollary 4.4.
(a) First, consider any element $`A𝚄(𝒯_{}(𝒦_{}(𝔖))^{\mathrm{po}})`$. According to Proposition 6.2, its symbol $`a`$ is (stably) homotopic to the constant loop $`1`$. Applying Proposition 6.5 to this homotopy, we see that it is sufficient to prove that Toeplitz units with symbol $`𝖴(a)\mathrm{\Lambda }(𝗓,𝖰)𝖴(a)^1\mathrm{\Lambda }(𝗓,𝖰)^1`$ can be contracted.
(b) Consider, again, $`A`$ as above. Let
$$Q=\left[\begin{array}{cc}1_{}& \\ & 𝖰\end{array}\right]=\left[\begin{array}{cccc}1_{}& & & \\ & 1_{}& & \\ & & & \\ & & & \\ & & & \\ & & 1_{}& \\ & & & 1_{}\end{array}\right];$$
here the double lines show how we decompose this block matrix of $`\times `$ matrices to a block matrix of $`\times `$ matrices. Furthermore, let
(8)
$$\begin{array}{c}S(\theta )=\left[\begin{array}{cccc}s& & & t\\ & 1& & \\ & & 1& \\ t& & & s\end{array}\right]\left[\begin{array}{ccccc}1_{}& & & & \\ & 1_{}& & \\ & & 𝖶(a^{})& 𝖸(a^{})\\ & & 𝖸(a)& 𝖶(a)\end{array}\right]\hfill \\ \hfill \left[\begin{array}{ccccc}1_{}& & & & \\ & 𝖶(a^1)& 𝖸(a^1)& \\ & 𝖸((a^1)^{})& 𝖶((a^1)^{})& \\ & & & 1_{}\end{array}\right]\left[\begin{array}{cccc}1_{}& & & \\ & A& & \\ & & 1_{}& \\ & & & A^1\end{array}\right]\left[\begin{array}{cccc}s& & & t\\ & 1& & \\ & & 1& \\ t& & & s\end{array}\right]\\ \hfill 𝚄(𝒦_{\{1,2\}\times }(𝔖)),\end{array}$$
and take $`S(\theta )L(Q,S(\theta )^1)𝚄(𝒯_{}(𝒦_{\{1,2\}\times }(𝔖)))`$. This yields a homotopy between $`S(0)L(Q,S(0)^1)`$ and $`S(\pi /2)L(Q,S(\pi /2)^1)`$, which have symbols
$$S(0)\mathrm{\Lambda }(𝗓,Q)S(0)^1\mathrm{\Lambda }(𝗓,Q)^1=\left[\begin{array}{cc}1_{}& \\ & 𝖴(a)\mathrm{\Lambda }(𝗓,𝖰)𝖴(a)^1\mathrm{\Lambda }(𝗓,𝖰)^1\end{array}\right]$$
and
$$S(\pi /2)\mathrm{\Lambda }(𝗓,Q)S(\pi /2)^1\mathrm{\Lambda }(𝗓,Q)^1=\left[\begin{array}{cc}1_{}& \\ & 1_{}\end{array}\right],$$
respectively. Thus, Toeplitz units with symbol $`\mathrm{\Lambda }(𝗓^1,𝖰)𝖴(a)\mathrm{\Lambda }(𝗓,𝖰)𝖴(a)^1`$ can be deformed to Toeplitz units with trivial symbols. According to part (a), it is sufficient to show that elements with trivial symbol can be contracted.
(c) Now suppose that the symbol of a Toeplitz unit $`A`$ is $`1`$. According to standard stabilization arguments, we can assume that $`A=𝖤_{}(\stackrel{~}{k})`$, where $`\stackrel{~}{k}=\left[\begin{array}{cc}k_0& \\ & 1_{}\end{array}\right]𝚄(𝒦_{}(𝔖))`$. Let $`\stackrel{~}{k}(\theta )=\left[\begin{array}{cc}s& t\\ t& s\end{array}\right]\left[\begin{array}{cc}\stackrel{~}{k}_0& \\ & 1_{}\end{array}\right]\left[\begin{array}{cc}s& t\\ t& s\end{array}\right].`$ Then $`\stackrel{~}{k}(\theta )L(𝖰,\stackrel{~}{k}(\theta ))^1`$ yields a homotopy between $`\stackrel{~}{k}(0)L(𝖰,\stackrel{~}{k}(0))^1=𝖤_{}(\stackrel{~}{k})=A`$ and $`\stackrel{~}{k}(\pi /2)L(𝖰,\stackrel{~}{k}(\pi /2))^1=1`$. ∎
###### Remark 6.7.
(a) If the locally convex algebra $`𝔄`$ is strong in the terminology of in , i. e. for all seminorm $`p`$ there is a seminorm $`\stackrel{~}{p}`$ such that $`p(X_1\mathrm{}X_n)\stackrel{~}{p}(X_1)\mathrm{}\stackrel{~}{p}(X_n)`$ holds for all $`n`$, then the proof can be much simplified: In that case, the associated algebras are also strong, and the smooth homotopy lifting property holds for the symbol map. Then, using Proposition 6.2, the proof of the contractibility statement reduces to point (c) immediately, hence making points (a) and (b), and the construction of 6.4 unnecessary. One must note that Proposition 6.6 above is much easier to prove than Kuiper’s Theorem about the contractibility of the unitary group. See, e. g. .
(b) Stabilization was an important assumption in the previous statement. For example, $`𝚄(𝒯_{}()^{\mathrm{po}})`$ is not contractible, as it allows an extended, multiplicate determinant.
## 7. Possible modifications
Due to the nice properties of $`U(a,\theta ,𝗏)`$, Statement 1.4 can be seen in a rather straightforward manner. We remark that another such category is the category of Hilbert-Schmidt operators, used by Pressley and Segal, , Ch. 6. Furthermore, with some extra work, the transformation parameter $`\theta `$ (i. e. $`s`$ and $`t`$ jointly) can be replaced by $`t`$ entirely, extending the constructions as formal homotopies.
## 8. Algebraically finite cyclic loops
A practical disadvantage of $`𝖡(a)`$ is that it is, in general, an infinite perturbation of $`𝖰`$. The exception is when $`a𝚄(𝔄[𝗓^1,𝗓]^\mathrm{f})`$, but this is a rather restrictive condition from geometrical viewpoint. We will show below that we can do well also in the case when $`a`$ can be represented by finite loops but it is not in $`𝚄(𝔄[𝗓^1,𝗓]^\mathrm{f})`$.
###### 8.1.
For $`m0n`$, we say that the loop $`a(𝗓)𝚄(𝔄[𝗓^1,𝗓])`$ is an $`L(m,n)`$-finite loop if $`a(𝗓)=_{mjn}a_j𝗓^k`$. A loop $`a(𝗓)`$ is an $`R(m,n)`$-finite loop if its inverse $`a(𝗓)^1`$ is an $`L(n,m)`$-finite loop. For a finite sequence $`F=\{(m_j,n_j)\}_{1js}`$, let
$$𝔄_F=\{(a_s,\mathrm{},a_1):a_j𝚄(𝔄[𝗓^1,𝗓])\text{ is }L(m_j,n_j)\text{ or }R(m_j,n_j)\text{-finite}\}.$$
We say that $`a𝚄(𝔄[𝗓^1,𝗓]`$ is algebraically finite of type $`F`$ if $`a=a_s\mathrm{}a_1`$ for an element $`(a_s,\mathrm{},a_1)𝔄_F`$.
###### 8.2.
For $`m0n`$, we say that a matrix $`A`$ is an $`L(m,n)`$-perturbation of $`A_0`$ if
$$A=A_0+\underset{min,j}{}a_{i,j}𝐞_{i,j},$$
for $`a_{i,j}`$ chosen suitably. Similarly, we can define $`R(m,n)`$-perturbations by interchanging the role of $`i`$ and $`j`$ in the expression above. An $`(m,n)`$-perturbation is a matrix which is both an $`L(m,n)`$-perturbation and an $`R(m,n)`$-perturbation.
In what follows, we will always be concerned with perturbations of $`\mathrm{\Lambda }(𝗌,𝖰)`$, where $`𝗌`$ is equal to $`1`$, $`1`$, or another formal variable $`𝗏`$. Both $`L(m,n)`$-perturbations and $`R(m,n)`$-perturbations of $`\mathrm{\Lambda }(𝗌,𝖰)`$ can be reduced to $`(m,n)`$-perturbations by taking direct cut-offs of unwanted matrix elements:
$$\left[\begin{array}{ccc}𝗌& & \\ L^{}& M& L^+\\ & & 1\end{array}\right]\stackrel{\mathrm{R}_{(m,n)}}{}\left[\begin{array}{ccc}𝗌& & \\ & M& \\ & & 1\end{array}\right]\stackrel{\mathrm{R}_{(m,n)}}{}\left[\begin{array}{ccc}𝗌& R^{}& \\ & M& \\ & R^+& 1\end{array}\right].$$
The reduction $`\mathrm{R}_{(m,n)}`$ is essentially taking away the off-diagonal elements of a triangular block matrix (with respect to an appropriate ordering of the basis). Sometimes it is practical to use the partial reduction $`\mathrm{R}_{(m,n)}^{[h]}=(1h)\mathrm{Id}+h\mathrm{R}_{(m,n)}`$, where $`h`$ is assumed to be a scalar variable. Here the off-diagonal blocks are not taken away completely but multiplied by $`1h`$. It is useful to notice that (partial) reduction is a homomorphism as long as we restrict our attention to matrices of appropriate block triangular shape. In particular, invertible elements / involutions are reduced to invertible elements / involutions.
###### 8.3.
The involutions $`Q`$ and $`\overline{Q}`$ are unipotently related if $`\frac{1}{2}(Q\overline{Q}+\overline{Q}Q)=1`$ holds. In this case the expression $`C(\overline{Q},Q)=\frac{1+\overline{Q}Q}{2}`$ satisfies the identities
$$C(\overline{Q},Q)^1=C(Q,\overline{Q})\text{and}C(\overline{Q},Q)QC(\overline{Q},Q)^1=\overline{Q}.$$
More generally, $`C(\overline{Q},Q,h)=(1h)1+h\frac{1+\overline{Q}Q}{2}`$ satisfies the identities
$$C(\overline{Q},Q,h)^1=C(Q,\overline{Q},h)\text{and}C(\overline{Q},Q,h)QC(\overline{Q},Q,h)^1=(1h)Q+h\overline{Q}.$$
This situation applies when, in the manner of the previous paragraph, an involution $`Q`$ is reduced to an involution $`\overline{Q}`$.
###### Lemma 8.4.
If $`a(𝗓)=_{mjn}a_j𝗓^k`$, $`m0n`$, then $`U(a,\theta ,𝗏)`$ is an $`(m,n)`$-perturbation of $`𝖴(a)`$.
###### Proof.
This is immediate from 5.2. ∎
###### Lemma 8.5.
Suppose that $`A`$ is an $`(m^{},n^{})`$-perturbation of $`\mathrm{\Lambda }(𝗌,𝖰)`$, where $`m^{}0n^{}`$. Then we claim:
If $`a`$ is an $`L(m,n)`$\- or $`R(m,n)`$-finite loop, then $`U(a,\theta _1,𝗏)AU(a,\theta _2,𝗐)^1`$ is an $`L(m+m^{},n+n^{})`$\- or $`R(m+m^{},n+n^{})`$-perturbation of $`\mathrm{\Lambda }(𝗌,𝖰)`$, respectively.
###### Proof.
The $`L`$ case: Let $`k>n+n^{}`$ and $`h=𝗌`$ if $`k<m+m^{}`$. The special shape of the matrices implies
$$𝐞_k^{}U(a,\theta _1,𝗏)A=\left(\underset{mjn}{}a_j𝐞_{kj}^{}\right)A=h\underset{mjn}{}a_j𝐞_{kj}^{}=h𝐞_k^{}U(a,\theta _2,𝗐),$$
from which $`𝐞_k^{}U(a,\theta _1,𝗏)AU(a,\theta _2,𝗐)^1=h𝐞_k^{}=𝐞_k^{}\mathrm{\Lambda }(𝗌,𝖰)`$. This latter equality, which holds for appropriate $`k`$, is exactly the statement of having an $`L(m+m^{},n+n^{})`$-perturbation of $`\mathrm{\Lambda }(𝗌,𝖰)`$. The $`R`$ case is similar. ∎
###### 8.6.
Next, we construct a linearization procedure which linearizes algebraically finite loops into finite perturbations: Let $`F=\{(m_j,n_j)\}_{1js}`$ be a finiteness type, $`\stackrel{~}{a}=(a_s,\mathrm{},a_1)𝔄_F`$, and $`a=a_s\mathrm{}a_1`$. Set $`M_k=m_1+\mathrm{}+m_k`$, $`N_k=n_1+\mathrm{}+n_k`$. Let $`|F|=(M_s,N_s)`$. Also, let $`\stackrel{~}{a}_k=(a_k,\mathrm{},a_1)`$, with appropriate finiteness type $`F_k`$. Then $`|F_k|=(M_k,N_k)`$. We define
$$𝖡_F(\stackrel{~}{a})=\mathrm{R}_{|F_s|}\left(𝖴(a_s)\mathrm{}\mathrm{R}_{|F_1|}\left(𝖴(a_1)\mathrm{𝖰𝖴}(a_1)^1\right)\mathrm{}𝖴(a_s)^1\right).$$
Then $`𝖡_F(\stackrel{~}{a})`$ is an involution, and an $`|F|`$-perturbation of $`𝖰`$. More generally, let
(9)
$$\begin{array}{c}K_F^\mathrm{c}(\stackrel{~}{a},\theta ,𝗏,𝗐)=\mathrm{R}_{|F_s|}(U(a_s,\theta ,𝗏)\mathrm{}\hfill \\ \hfill \mathrm{}\mathrm{R}_{|F_1|}\left(U(a_1,\theta ,𝗏)\mathrm{\Lambda }(\mathrm{𝗏𝗐}^1,𝖰)U(a_1,\theta ,𝗐)^1\right)\mathrm{}U(a_s,\theta ,𝗐)^1).\end{array}$$
Then, in particular, $`K_F^\mathrm{c}(\stackrel{~}{a},0,𝗏,𝗐)=\mathrm{\Lambda }(\mathrm{𝗏𝗐}^1,𝖡_F(\stackrel{~}{a}))`$, and $`K_F^\mathrm{c}(\stackrel{~}{a},\pi /2,𝗏,𝗐)=𝖤_{}(a(𝗓))\mathrm{\Lambda }(\mathrm{𝗏𝗐}^1,𝖰)𝖤_{}(a(𝗐)^1)`$; which are immediate from the special shape of the matrices involved. This yields
###### Proposition 8.7.
The continuous map
$$K_F:𝔄_F\times \mathrm{S}^1𝚄(𝒦_{}(𝔄)[𝗏^1,𝗏]^{\mathrm{po}})$$
defined by
$$\stackrel{~}{a},\theta K_F(\stackrel{~}{a},\theta ,𝗏)=K_F^\mathrm{c}(\stackrel{~}{a},\theta ,𝗏,1)\mathrm{\Lambda }(𝗏,𝖰)^1$$
is smooth in the variable $`\theta `$; and it is an $`|F|`$-perturbation of $`1_{}`$. Here
$$K_F(\stackrel{~}{a},0,𝗏)=\mathrm{\Lambda }(𝗏,𝖡_F(\stackrel{~}{a}))\mathrm{\Lambda }(𝗏,𝖰)^1,K_F(\stackrel{~}{a},\pi /2,𝗏)=𝖤_{}(a(𝗏)a(1)^1).$$
In particular, as $`\mathrm{S}^1`$ is restricted to $`[0,\pi /2]`$, it yields a linearizing homotopy of $`a(𝗓)a(1)^1`$ in the finite perturbation category. ∎
In the literature one finds comments about the possibly very large size of the matrices used in linearizing homotopies. The result above, however, shows the one can do reasonably well.
###### 8.8.
There is, however, a closer analogy between the non-finite and the finite cases: Let $`Q_0=𝖰`$, and $`Q_k=\mathrm{R}_{|F_k|}(𝖴(a_k)Q_{k1}𝖴(a_k)^1)`$ by recursion. Then $`Q_k=𝖡_{F_k}(\stackrel{~}{a}_k)`$. Using the notation $`_{i=1}^sx_i=x_n\mathrm{}x_2x_1`$, let
$$𝖴_F(\stackrel{~}{a})=\underset{i=1}{\overset{s}{}}\frac{𝖴(a_i)+Q_i𝖴(a_i)Q_{i1}}{2}=\underset{i=1}{\overset{s}{}}C(Q_i,𝖴(a_i)Q_{i1}𝖴(a_i)^1)𝖴(a_i).$$
According to our earlier observations,
$$𝖡_F(\stackrel{~}{a})=𝖴_F(\stackrel{~}{a})\mathrm{𝖰𝖴}_F(\stackrel{~}{a})^1.$$
We also define
$$U_F(\stackrel{~}{a},\theta ,𝗏)=\mathrm{R}_{|F_s|}\left(U(a_s,\theta ,𝗏)\mathrm{}\mathrm{R}_{|F_1|}\left(U(a_1,\theta ,𝗏)𝖴(a_1)^1\right)\mathrm{}𝖴(a_s)^1\right)𝖴_F(\stackrel{~}{a}).$$
and
$$\widehat{𝖴}_F(\stackrel{~}{a})=\mathrm{R}_{|F_s|}\left(\widehat{𝖴}(a_s)\mathrm{}\mathrm{R}_{|F_1|}\left(\widehat{𝖴}(a_1)𝖴(a_1)^1\right)\mathrm{}𝖴(a_s)^1\right)𝖴_F(\stackrel{~}{a}).$$
Then $`U_F(\stackrel{~}{a},0,𝗏)=𝖴_F(\stackrel{~}{a}),`$ which is trivial; and, analogously to the original situation, $`U_F(\stackrel{~}{a},\pi /2,𝗏)=𝖤_{}(a(𝗓))\mathrm{\Lambda }(𝗏,𝖰)\widehat{𝖴}_F(\stackrel{~}{a})\mathrm{\Lambda }(𝗏,𝖰)^1,`$ which follows from $`\mathrm{\Lambda }(𝗏,𝖡_F(\stackrel{~}{a}))^1=𝖴_F(\stackrel{~}{a})\mathrm{\Lambda }(𝗏,𝖰)^1𝖴_F(\stackrel{~}{a})^1`$ and the homomorphism property of reduction. In fact,
$$K_F^\mathrm{c}(\stackrel{~}{a},\theta ,𝗏,𝗐)=U_F(\stackrel{~}{a},\theta ,𝗏)\mathrm{\Lambda }(\mathrm{𝗏𝗐}^1,𝖰)U_F(\stackrel{~}{a},\theta ,𝗐)^1$$
holds. Again, this follows from $`\mathrm{\Lambda }(\mathrm{𝗏𝗐}^1,𝖡_F(\stackrel{~}{a}))=𝖴_F(\stackrel{~}{a})\mathrm{\Lambda }(\mathrm{𝗏𝗐}^1,𝖰)𝖴_F(\stackrel{~}{a})^1`$ and the homomorphism property of reduction.
###### 8.9.
The constructions above can be expounded in order to show that the linearizations $`K`$ and $`K_F`$ can nicely be deformed into each other: Let
$$𝖴_F(\stackrel{~}{a},h)=\underset{k=1}{\overset{s}{}}(1h)𝖴(a_k)+h𝖴_{F_k}(\stackrel{~}{a}_k)𝖴_{F_{k1}}(\stackrel{~}{a}_{k1})^1$$
Here the product terms can also be written as $`C(Q_k,𝖴(a_k)Q_{k1}𝖴(a_k)^1,h)𝖴(a_k)`$, which makes invertibility clear. Then $`𝖴_F(\stackrel{~}{a},0)=𝖴(a)`$, $`𝖴_F(\stackrel{~}{a},1)=𝖴_F(\stackrel{~}{a})`$. Let
$$𝖡_F(\stackrel{~}{a},h)=𝖴_F(\stackrel{~}{a},h)\mathrm{𝖰𝖴}_F(\stackrel{~}{a},h)^1.$$
Notice that $`𝖡_F(\stackrel{~}{a},0)=𝖡(a)`$, $`𝖡_F(\stackrel{~}{a},1)=𝖡_F(\stackrel{~}{a}).`$ Let
(10)
$$\begin{array}{c}U_F(\stackrel{~}{a},h,\theta ,𝗏)=\underset{k=1}{\overset{s}{}}\left((1h)U(a_k,\theta ,𝗏)+hU_{F_k}(\stackrel{~}{a}_k,\theta ,𝗏)U_{F_{k1}}(\stackrel{~}{a}_{k1},\theta ,𝗏)^1\right)=\hfill \\ \hfill =\underset{k=1}{\overset{s}{}}\mathrm{R}_{|F_k|}^{[h]}\left(U(a_k,\theta ,𝗏)\mathrm{}\mathrm{R}_{|F_1|}\left(U(a_1,\theta ,𝗏)𝖴(a_1)^1\right)\mathrm{}𝖴(a_k)^1\right)\\ \hfill C(Q_k,𝖴(a_k)Q_{k1}𝖴(a_k)^1,h)𝖴(a_k)\\ \hfill \mathrm{R}_{|F_{k1}|}\left(U(a_{k1},\theta ,𝗏)\mathrm{}\mathrm{R}_{|F_1|}\left(U(a_1,\theta ,𝗏)𝖴(a_1)^1\right)\mathrm{}𝖴(a_{k1})^1\right)^1.\end{array}$$
Again, the latter product form implies not only invertibility but that the inverses of the product terms are linear in $`h`$. In particular, it yields that the inverse is
$$U_F(\stackrel{~}{a},h,\theta ,𝗏)^1=\underset{k=1}{\overset{s}{}}(1h)U(a_k,\theta ,𝗏)^1+hU_{F_{k1}}(\stackrel{~}{a}_{k1},\theta ,𝗏)U_{F_k}(\stackrel{~}{a}_k,\theta ,𝗏)^1.$$
This also shows that $`U_F(\stackrel{~}{a},h,\theta ,𝗏)^1`$ is polynomial in $`h`$. We also define
$$\widehat{𝖴}_F(\stackrel{~}{a},h)=\underset{k=1}{\overset{s}{}}\left((1h)\widehat{𝖴}(a_k)+h\widehat{𝖴}_{F_k}(\stackrel{~}{a}_k)\widehat{𝖴}_{F_{k1}}(\stackrel{~}{a}_{k1})^1\right)$$
One can see that the identities $`U_F(\stackrel{~}{a},h,0,𝗏)=𝖴_F(\stackrel{~}{a},h)`$ and
$$U_F(\stackrel{~}{a},h,\pi /2,𝗏)=𝖤_{}(a(𝗓))\mathrm{\Lambda }(𝗏,𝖰)\widehat{𝖴}_F(\stackrel{~}{a},h)\mathrm{\Lambda }(𝗏,𝖰)^1$$
hold. Furthermore, $`U_F(\stackrel{~}{a},0,\theta ,𝗏)=U(a,\theta ,𝗏)`$, $`U_F(\stackrel{~}{a},1,\theta ,𝗏)=U_F(\stackrel{~}{a},\theta ,𝗏)`$, and $`\widehat{𝖴}_F(\stackrel{~}{a},0)=\widehat{𝖴}(a)`$, $`\widehat{𝖴}_F(\stackrel{~}{a},1)=\widehat{𝖴}_F(\stackrel{~}{a})`$. We define
$$K_F^{\mathrm{ec}}(\stackrel{~}{a},h,\theta ,𝗏,𝗐)=U_F(\stackrel{~}{a},h,\theta ,𝗏)\mathrm{\Lambda }(\mathrm{𝗏𝗐}^1,𝖰)U_F(\stackrel{~}{a},h,\theta ,𝗐)^1.$$
From the earlier observations, the identities $`K_F^{\mathrm{ec}}(\stackrel{~}{a},h,0,𝗏,𝗐)=\mathrm{\Lambda }(\mathrm{𝗏𝗐}^1,𝖡_F(\stackrel{~}{a},h))`$ and
$$K_F^{\mathrm{ec}}(\stackrel{~}{a},h,\pi /2,𝗏,𝗐)=𝖤_{}(a(𝗏))\mathrm{\Lambda }(\mathrm{𝗏𝗐}^1,𝖰)𝖤_{}(a(𝗐))^1$$
follow. Furthermore, $`K_F^{\mathrm{ec}}(\stackrel{~}{a},0,\theta ,𝗏,𝗐)=K^\mathrm{c}(a,\theta ,𝗏,𝗐)`$ and $`K_F^{\mathrm{ec}}(\stackrel{~}{a},1,\theta ,𝗏,𝗐)=K_F^\mathrm{c}(\stackrel{~}{a},\theta ,𝗏,𝗐)`$.
This yields
###### Proposition 8.10 ($``$ Statement 1.5).
The continuous map
$$K_F^\mathrm{e}:𝔄_F\times \times \mathrm{S}^1𝚄(𝒦_{}(𝔄)[𝗏^1,𝗏]^{\mathrm{po}})$$
defined by
$$\stackrel{~}{a},h,\theta K_F^\mathrm{e}(\stackrel{~}{a},h,\theta ,𝗏)=K_F^{\mathrm{ec}}(\stackrel{~}{a},h,\theta ,𝗏,1)\mathrm{\Lambda }(𝗏,𝖰)^1$$
is smooth in $`\theta `$ and polynomial in $`h`$. It has the properties
(i) $`K_F^\mathrm{e}(\stackrel{~}{a},0,\theta )=K(a,\theta )`$;
(ii) $`K_F^\mathrm{e}(\stackrel{~}{a},1,\theta )=K_F(\stackrel{~}{a},\theta )`$;
(iii) $`K_F^\mathrm{e}(\stackrel{~}{a},h,0)=\mathrm{\Lambda }(𝗏,𝖡_F(\stackrel{~}{a},h)\mathrm{\Lambda }(𝗏,𝖰)^1`$;
(iv) $`K_F^\mathrm{e}(\stackrel{~}{a},h,\pi /2)=𝖤_{}(a(𝗓)a(1)^1)`$.
In particular, it connects the pullback homotopy $`K|_{𝔄_F}`$ and homotopy $`K_F`$ through other linearizing homotopies. ∎
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# Optimal Multiple Assignments Based on Integer Programming in Secret Sharing Schemes with General Access StructuresThis work has been submitted to the IEEE for possible publication. Copyright may be transferred without notice after which this version may no longer be accessible.
## 1 Introduction
A Secret Sharing Scheme (SSS) is a method to encrypt a secret information $`S`$ into $`n`$ pieces called shares $`V_1,V_2,\mathrm{},V_n`$, each of which has no information of the secret $`S`$, but $`S`$ can be decrypted by collecting several shares. For example, a $`(k,n)`$-threshold SSS means that any $`k`$ out of $`n`$ shares can decrypt the secret $`S`$ although any $`k1`$ or less shares do not leak out any information of $`S`$. The $`(k,n)`$-threshold access structure can be generalized to so-called general access structures which consist of the families of qualified sets and forbidden sets. A qualified set is the subset of shares that can decrypt the secret, but a forbidden set is the subset that does not leak out any information of $`S`$.
Generally, the efficiency of a SSS is measured by the entropy of each share. It is known that for any access structures, the entropies of secret $`S`$ and shares $`V_i`$, $`i=1,2,\mathrm{},n`$, must satisfy $`H(V_i)H(S)`$ . On the other hand, in the case of $`(k,n)`$-threshold SSSs, the optimal SSSs attaining $`H(V_i)=H(S)`$ can easily be constructed . However, it is hard to derive efficient SSSs for arbitrarily given general access structures although several construction methods have been proposed.
For example, the monotone circuit construction is a method to realize a SSS by combining several $`(m,m)`$-threshold SSSs. This method is simple but inefficient, and hence, it is extended to the decomposition construction , which uses several decomposed general SSSs. Although the decomposition construction can attain the optimal coding rates for some special access structures, it cannot construct an efficient SSS in the case that the decomposed SSSs cannot be realized efficiently. Note that a monotone circuit construction is based on qualified sets. Hence, as another extension of monotone circuit construction, a method is proposed to construct a SSS with general access structures based on qualified sets and $`(t,m)`$-threshold SSSs .
On the other hand, for any given general access structure, a SSS can be constructed from a $`(t,m)`$-threshold SSS by a multiple assignment map such that $`t`$ or more shares of the $`(t,m)`$-threshold SSS are assigned to qualified sets but $`t1`$ or less shares are assigned to forbidden sets. The cumulative map is a simple realization of the multiple assignment map based on an $`(m,m)`$-threshold SSS , and from the simplicity, it is often used in visual secret sharing schemes for general access structures . However, it is known that the SSS constructed by the cumulative map is inefficient generally, especially in the case that the access structure is a $`(k,n)`$-threshold SSS with $`kn`$. Recently, a modified cumulative map based on a $`(t,m)`$-threshold SSS is proposed to overcome this defect . But, the modified cumulative map is not always more efficient than the original cumulative map.
In this paper, we propose a new construction method that can derive the optimal multiple assignment map by integer programming. The proposed construction method is simple and optimal in the sense of multiple assignment maps. Furthermore, it can also be applied to incomplete and/or ramp access structures.
This paper is organized as follows. In Section 2, we give the definitions of SSSs and introduce the multiple assignment map. We also introduce the construction methods of the cumulative map and the modified cumulative map, and we point out their defects. To overcome such defects, we propose a new construction method of the optimal multiple assignment map by integer programming in Section 3. Finally, Sections 4 and 5 are devoted to present the applications of the proposed method to incomplete or ramp SSSs for general access structures, respectively.
## 2 Preliminaries
### 2.1 Definitions
Throughout this paper, a set of shares and a family of share sets are represented by bold-face and script letters, respectively. For sets $`𝑨`$ and $`𝑩`$, we denote a difference set by $`𝑨𝑩`$, which is defined as $`𝑨𝑩\stackrel{\text{def}}{=}𝑨\overline{𝑩}`$ where $`\overline{𝑩}`$ means the complement of $`𝑩`$. Furthermore, the cardinality of $`𝑨`$ is represented by $`|𝑨|`$, and the Cartesian product of $`𝑨`$ and $`𝑩`$ is expressed by $`𝑨\times 𝑩`$.
Let $`𝑽=\{V_1,V_2,\mathrm{},V_n\}`$ be the set of shares, and let $`2^𝑽`$ be the family of all subsets of $`𝑽`$. We represent the family of qualified sets that can decrypt a secret information $`S`$ and the family of forbidden sets that cannot gain any information of $`S`$ by $`𝒜_1`$ and $`𝒜_0`$, respectively.
$`\mathrm{\Gamma }=\{𝒜_1,𝒜_0\}`$ is called an access structure. For instance, the access structure of $`(k,n)`$-threshold SSSs can be represented as follows:
$`𝒜_1`$ $`=`$ $`\{𝑨2^𝑽:k|𝑨|n\},`$ (1)
$`𝒜_0`$ $`=`$ $`\{𝑨2^𝑽:0|𝑨|k1\}.`$ (2)
In SSSs, it obviously holds that $`𝒜_1𝒜_0=\mathrm{}`$. If it also holds that $`𝒜_1𝒜_0=2^𝑽`$, the access structure is called complete. Note that any access structure must satisfy the following monotonicity.
$`𝑨𝒜_1𝑨^{}𝒜_1\text{for all}𝑨^{}𝑨`$ (3)
$`𝑨𝒜_0𝑨^{}𝒜_0\text{for all}𝑨^{}𝑨`$ (4)
Therefore, we can define the family of minimal qualified sets and the family of maximal forbidden sets as follows:
$`𝒜_1^{}`$ $`=`$ $`\{𝑨𝒜_1:𝑨\{V\}𝒜_1\text{for any}V𝑨\},`$ (5)
$`𝒜_0^+`$ $`=`$ $`\{𝑨𝒜_0:𝑨\{V\}𝒜_0\text{for any}V𝑽𝑨\}.`$ (6)
We assume that the secret information $`S`$ and each share $`V_i`$ are random variables, which take values in finite fields $`𝔽_S`$ and $`𝔽_{V_i}`$, respectively. Then, share set $`𝑨=\{V_{i_1},V_{i_2},\mathrm{},V_{i_u}\}(𝑽)`$, which takes values in $`𝔽_𝑨\stackrel{\text{def}}{=}𝔽_{V_{i_1}}\times 𝔽_{V_{i_2}}\times \mathrm{}\times 𝔽_{V_{i_u}}`$, must satisfy the following conditions:
$`H(S|𝑨)`$ $`=`$ $`H(S)\text{if}𝑨𝒜_0,`$ (7)
$`H(S|𝑨)`$ $`=`$ $`0\text{if}𝑨𝒜_1,`$ (8)
where $`H(S)`$ is the entropy of $`S`$ and $`H(S|𝑨)`$ is the conditional entropy of $`S`$ for given $`𝑨`$.
Now, let us define the coding rate of a share $`V_i`$ as $`\rho _i\stackrel{\text{def}}{=}H(V_i)/H(S)`$, for $`i=1,2,\mathrm{},n`$. Since each $`\rho _i`$ may be different in the case of general access structures, it is cumbersome to treat each $`\rho _i`$ independently. Hence, we consider only the following average coding rate $`\stackrel{~}{\rho }`$ and worst coding rate $`\rho ^{}`$.
$`\stackrel{~}{\rho }`$ $`\stackrel{\text{def}}{=}`$ $`{\displaystyle \frac{1}{n}}{\displaystyle \underset{i=1}{\overset{n}{}}}\rho _i,`$ (9)
$`\rho ^{}`$ $`\stackrel{\text{def}}{=}`$ $`\underset{1in}{\mathrm{max}}\rho _i.`$ (10)
For a given access structure $`\mathrm{\Gamma }=\{𝒜_1,𝒜_0\}`$, we call $`V𝑽`$ a significant share if there exists a share set $`𝑨2^𝑽`$ such that $`𝑨\{V\}𝒜_1`$ but $`𝑨𝒜_0`$.
###### Remark 1
Note that a non-significant share plays no roll in the SSS, and hence, $`\rho _i=0`$ can always be attained for each non-significant share $`V_i`$ in any access structure $`\mathrm{\Gamma }`$. Furthermore, if there exists a non-significant share $`V_i`$ with $`\rho _i>0`$, the average coding rate can be reduced by setting $`\rho _i=0`$ without changing all the significant shares. Hence, we call a non-significant share a vacuous share. On the other hand, we have $`\rho _i1`$ for any significant share $`V_i`$ because it must satisfy $`H(V_i)H(S)`$ . In the following, we assume that every share is significant. $`\mathrm{}`$
If a SSS attains $`\rho _i=1`$ for all $`i`$, it is called ideal. It is known that in the case of $`(k,n)`$-threshold SSSs, the ideal SSS can easily be constructed for any $`k`$ and $`n`$ . Since $`\rho _i1`$, $`i=1,2,\mathrm{},n`$, must hold for any significant share $`V_i`$ in any access structures, $`\stackrel{~}{\rho }=1`$ or $`\rho ^{}=1`$ are the necessary and sufficient conditions for a SSS to be ideal .
### 2.2 Multiple Assignment Map
Let $`\mathrm{\Gamma }=\{𝒜_1,𝒜_0\}`$ be a given general access structure with share set $`𝑽=\{V_1,V_2,\mathrm{},V_n\}`$ and let $`𝑾_{(t,m)}=\{W_1^{(t)},W_2^{(t)},\mathrm{},W_m^{(t)}\}`$ be the share set of a $`(t,m)`$-threshold SSS. We now consider a map $`\phi _\mathrm{\Gamma }:\{1,2,\mathrm{},n\}2^{𝑾_{(t,m)}}`$, which assigns each participant a subset of the shares generated by the $`(t,m)`$-threshold scheme, and a map $`\mathrm{\Phi }_\mathrm{\Gamma }:2^𝑽2^{𝑾_{(t,m)}}`$, which is defined as $`\mathrm{\Phi }_\mathrm{\Gamma }(𝑨)\stackrel{\text{def}}{=}_{V_i𝑨}\phi _\mathrm{\Gamma }(i)`$ for a share set $`𝑨𝑽`$. Then, $`\phi _\mathrm{\Gamma }`$ is called a multiple assignment map for the access structure $`\mathrm{\Gamma }`$ if each share $`V_i`$ is determined by $`V_i=\phi _\mathrm{\Gamma }(i)`$ and $`\mathrm{\Phi }_\mathrm{\Gamma }(𝑨)`$ satisfies the following conditions:
$`|\mathrm{\Phi }_\mathrm{\Gamma }(𝑨)|`$ $``$ $`t\text{if}𝑨𝒜_1,`$ (11)
$`|\mathrm{\Phi }_\mathrm{\Gamma }(𝑨)|`$ $``$ $`t1\text{if}𝑨𝒜_0,`$ (12)
$`\mathrm{\Phi }_\mathrm{\Gamma }(𝑽)`$ $`=`$ $`𝑾_{(t,m)}.`$ (13)
To distinguish $`W_j^{(t)}𝑾_{(t,m)}`$ from the shares $`V_i`$ of $`\mathrm{\Gamma }`$, we call $`W_j^{(t)}`$ a primitive share.
Since any $`(t,m)`$-threshold SSS can easily be constructed as an ideal SSS , we assume in this paper that the $`(t,m)`$-threshold SSS with $`𝑾_{(t,m)}=\{W_1^{(t)},W_2^{(t)},\mathrm{},W_m^{(t)}\}`$ is ideal. Then, the average and worst coding rates defined by (9) and (10) become
$`\stackrel{~}{\rho }`$ $`=`$ $`{\displaystyle \frac{1}{n}}{\displaystyle \underset{i=1}{\overset{n}{}}}|\phi _\mathrm{\Gamma }(i)|,`$ (14)
$`\rho ^{}`$ $`=`$ $`\underset{1in}{\mathrm{max}}|\phi _\mathrm{\Gamma }(i)|,`$ (15)
respectively, since it holds that $`\rho _i=|\phi _\mathrm{\Gamma }(i)|`$.
In the case of $`t=m`$, it is known that the multiple assignment map $`\phi _\mathrm{\Gamma }`$ satisfying (11)–(13) can be realized for any access structures . Suppose that the access structure $`\mathrm{\Gamma }=\{𝒜_1,𝒜_0\}`$ has
$`𝒜_0^+=\{𝑭_1,𝑭_2,\mathrm{},𝑭_m\}.`$ (16)
Note that $`m=\left|𝒜_0^+\right|`$. Then, consider the map $`\psi _\mathrm{\Gamma }:\{1,2,\mathrm{},n\}2^{𝑾_{(m,m)}}`$ defined by
$`\psi _\mathrm{\Gamma }(i)={\displaystyle \underset{j:V_i𝑭_j}{}}\left\{W_j^{(m)}\right\}`$ (17)
where $`𝑭_j𝒜_0^+`$ and $`𝑾_{(m,m)}=\{W_1^{(m)},W_2^{(m)},\mathrm{},W_m^{(m)}\}`$ is the set of primitive shares of an $`(m,m)`$-threshold SSS. The above multiple assignment map $`\psi _\mathrm{\Gamma }`$ is called the cumulative map.
###### Example 2
Assume that $`n=4`$ and access structure $`\mathrm{\Gamma }_1`$ is defined by
$`𝒜_1^{}`$ $`=`$ $`\{\{V_1,V_2,V_3\},\{V_1,V_4\},\{V_2,V_4\},\{V_3,V_4\}\},`$ (18)
$`𝒜_0^+`$ $`=`$ $`\{\{V_1,V_2\},\{V_1,V_3\},\{V_2,V_3\},\{V_4\}\}.`$ (19)
Then, $`m=\left|𝒜_0^+\right|=4`$, and the cumulative map $`\psi _{\mathrm{\Gamma }_1}`$ is given from (17) as follows.
$`V_1`$ $`=`$ $`\psi _{\mathrm{\Gamma }_1}(1)=\{W_3^{(4)},W_4^{(4)}\},`$ (20)
$`V_2`$ $`=`$ $`\psi _{\mathrm{\Gamma }_1}(2)=\{W_2^{(4)},W_4^{(4)}\},`$ (21)
$`V_3`$ $`=`$ $`\psi _{\mathrm{\Gamma }_1}(3)=\{W_1^{(4)},W_4^{(4)}\},`$ (22)
$`V_4`$ $`=`$ $`\psi _{\mathrm{\Gamma }_1}(4)=\{W_1^{(4)},W_2^{(4)},W_3^{(4)}\}.`$ (23)
In this example, it holds that $`\stackrel{~}{\rho }=9/4`$ and $`\rho ^{}=3`$. $`\mathrm{}`$
It is known that the next theorem holds for the cumulative map $`\psi _\mathrm{\Gamma }`$.
###### Theorem 3 ()
For any multiple assignment map $`\phi _\mathrm{\Gamma }:\{1,2,\mathrm{},n\}2^{𝑾_{(t,m)}}`$ with $`t=m`$, it must hold that $`|𝑾_{(m,m)}||𝒜_0^+|`$, i.e., $`m|𝒜_0^+|`$. The equality holds if and only if $`\phi _\mathrm{\Gamma }(i)`$ is equal to the cumulative map $`\psi _\mathrm{\Gamma }(i)`$ defined by (17), where we assume that all $`\psi _\mathrm{\Gamma }`$’s obtained by permutations of $`𝑭_j`$’s in (16) are the same. $`\mathrm{}`$
Theorem 3 means that, in the case of $`t=m`$, the cumulative map $`\psi _\mathrm{\Gamma }`$ minimizes the number of primitive shares $`m`$. But, the minimization of $`m`$ does not mean the realization of an efficient SSS generally because it does not minimize the average coding rate $`\stackrel{~}{\rho }`$ and/or the worst coding rate $`\rho ^{}`$.
For instance, consider the case that $`\mathrm{\Gamma }`$ is a $`(k,n)`$-threshold access structure with $`kn`$. If we construct shares $`V_i`$ by the cumulative map $`\psi `$ for this $`\mathrm{\Gamma }`$, each $`V_i`$ must consist of $`\left(\genfrac{}{}{0pt}{}{n1}{k1}\right)`$ primitive shares of an $`(\left(\genfrac{}{}{0pt}{}{n}{k1}\right),\left(\genfrac{}{}{0pt}{}{n}{k1}\right))`$-threshold SSS because of $`|𝒜_0^+|=\left(\genfrac{}{}{0pt}{}{n}{k1}\right)`$. This means that $`\stackrel{~}{\rho }=\rho ^{}=\left(\genfrac{}{}{0pt}{}{n1}{k1}\right)`$. But, if we use the $`(k,n)`$-threshold SSS itself, we have $`\stackrel{~}{\rho }=\rho ^{}=1`$ because each $`V_i`$ consists of one primitive share. Hence, the cumulative map is quite inefficient in the case that $`\mathrm{\Gamma }`$ is a $`(k,n)`$-threshold access structure. In order to overcome this defect, a modified cumulative map is proposed in based on $`(t,m)`$-threshold SSSs. The modified cumulative map $`\psi _\mathrm{\Gamma }^{}`$ is constructed as follows.
###### Construction 4 ()
For a given $`\mathrm{\Gamma }=\{𝒜_0^+,𝒜_1^{}\}`$ and a positive integer $`g\stackrel{\text{def}}{=}\underset{𝑨𝒜_1^{}}{\mathrm{min}}|𝑨|`$, let $`𝒢_0𝒜_0^+`$ be the family defined by
$`𝒢_0=\{𝑮𝒜_0^+:|𝑮|g\}.`$ (24)
When $`𝒢_0=\{𝑮_1,𝑮_2,\mathrm{},𝑮_u\}\mathrm{}`$, let $`l_j\stackrel{\text{def}}{=}|𝑮_j|g+1`$ for $`j=1,2,\mathrm{},u`$, and $`\mathrm{}_j\stackrel{\text{def}}{=}_{p=1}^jl_p`$. If $`𝒢_0=\mathrm{}`$, let $`u=1`$ and $`\mathrm{}_1=0`$. Then, consider a $`(g+\mathrm{}_u,n+\mathrm{}_u)`$-threshold SSS and the set of primitive shares $`𝑾_{(g+\mathrm{}_u,n+\mathrm{}_u)}=\{W_1^{(g+\mathrm{}_u)},W_2^{(g+\mathrm{}_u)},\mathrm{},W_{n+\mathrm{}_u}^{(g+\mathrm{}_u)}\}`$. Furthermore, let $`𝑼_j`$, $`j=1,2,\mathrm{},u`$, be the subset of primitive shares defined by
$`𝑼_1`$ $`=`$ $`\mathrm{}\text{if}𝒢_0=\mathrm{},`$ (25)
$`𝑼_j`$ $`=`$ $`\{W_{n+\mathrm{}_{j1}+1}^{(g+\mathrm{}_u)},W_{n+\mathrm{}_{j1}+2}^{(g+\mathrm{}_u)},\mathrm{},W_{n+\mathrm{}_j}^{(g+\mathrm{}_u)}\}\text{if}𝒢_0\mathrm{},`$ (26)
where $`\mathrm{}_0=0`$. Then, the modified cumulative map $`\psi _\mathrm{\Gamma }^{}`$ is defined by
$`\psi _\mathrm{\Gamma }^{}(i)=\left\{W_i^{(g+\mathrm{}_u)}\right\}\left\{{\displaystyle \underset{j:V_i𝑮_j}{}}𝑼_j\right\}.`$ (27)
$`\mathrm{}`$
In the case where $`\mathrm{\Gamma }`$ is a $`(k,n)`$-threshold access structure, it holds that $`𝒢_0=\mathrm{}`$ and $`𝑼_1=\mathrm{}`$, and hence, it holds that $`\psi _\mathrm{\Gamma }^{}(i)=\{W_i^{(k)}\}`$ for $`i=1,2,\mathrm{},n`$ and this scheme coincides with the ideal $`(k,n)`$-threshold SSS . Therefore, the modified cumulative map $`\psi _\mathrm{\Gamma }^{}`$ is efficient if $`\mathrm{\Gamma }`$ is, or is near to, a $`(k,n)`$-threshold access structures. Furthermore, it is shown in that if the access structure $`\mathrm{\Gamma }`$ satisfies
$`\left|𝒜_0^+\right|{\displaystyle \frac{(ng1)\mathrm{}_u+n+2|𝒢_0|}{ng+1}},`$ (28)
then it holds that for the original cumulative map $`\psi _\mathrm{\Gamma }`$, $`_{V_i𝑽}|\psi _\mathrm{\Gamma }^{}(i)|_{V_i𝑽}|\psi _\mathrm{\Gamma }(i)|`$, which means that the average coding rate $`\stackrel{~}{\rho }`$ of $`\psi _\mathrm{\Gamma }^{}`$ is smaller than or equal to $`\psi _\mathrm{\Gamma }`$.
But, as shown in the following example, $`\psi _\mathrm{\Gamma }^{}`$ is not always more efficient than $`\psi _\mathrm{\Gamma }`$ if $`\mathrm{\Gamma }`$ does not satisfy (28).
###### Example 5
Consider the access structure $`\mathrm{\Gamma }_1`$ given by (18) and (19) in Example 2, which does not satisfy (28). Since we have $`g=2`$ from (18), $`𝒢_0`$ becomes $`𝒢_0=\{\{V_1,V_2\},\{V_1,V_3\},\{V_2,V_3\}\}\stackrel{\text{def}}{=}\{𝑮_1,𝑮_2,𝑮_3\}`$. Furthermore, since we have that $`l_1=l_2=l_3=1`$ and $`\mathrm{}_3=3`$, $`𝑼_i`$’s are determined as $`𝑼_1=\{W_5^{(5)}\},𝑼_2=\{W_6^{(5)}\},𝑼_3=\{W_7^{(5)}\}`$ for $`𝑾_{(5,7)}=\{W_1^{(5)},W_2^{(5)},\mathrm{},W_7^{(5)}\}`$. Hence, we can check that $`\mathrm{\Gamma }_1`$ does not satisfy (28) because of $`|𝒜_0^+|=4`$, $`n=4`$, $`g=2`$, $`\mathrm{}_u=3`$, and $`|𝒢_0|=3`$. Finally, we have from (27) that
$`V_1`$ $`=`$ $`\psi _{\mathrm{\Gamma }_1}^{}(1)=\{W_1^{(5)},W_7^{(5)}\},`$ (29)
$`V_2`$ $`=`$ $`\psi _{\mathrm{\Gamma }_1}^{}(2)=\{W_2^{(5)},W_6^{(5)}\},`$ (30)
$`V_3`$ $`=`$ $`\psi _{\mathrm{\Gamma }_1}^{}(3)=\{W_3^{(5)},W_5^{(5)}\},`$ (31)
$`V_4`$ $`=`$ $`\psi _{\mathrm{\Gamma }_1}^{}(4)=\{W_4^{(5)},W_5^{(5)},W_6^{(5)},W_7^{(5)}\}.`$ (32)
In this example, the coding rates are given by $`\stackrel{~}{\rho }=5/2`$ and $`\rho ^{}=4`$, which are larger than the coding rates of Example 2, i.e., $`\stackrel{~}{\rho }=9/4`$ and $`\rho ^{}=3`$. $`\mathrm{}`$
Note that (28) does not guarantee that the worst coding rate $`\rho ^{}`$ of $`\psi _\mathrm{\Gamma }^{}`$ is smaller than $`\psi _\mathrm{\Gamma }`$. Actually, the next example shows a case where $`\psi _\mathrm{\Gamma }^{}`$ attains a smaller average coding rate but gives larger worst coding rate than $`\psi _\mathrm{\Gamma }`$.
###### Example 6
Consider the access structure $`\mathrm{\Gamma }_2`$ given by
$`𝒜_1^{}`$ $`=`$ $`\{\{V_1,V_2,V_3,V_5\},\{V_1,V_2,V_4\},\{V_1,V_3,V_4\},\{V_1,V_4,V_5\},`$ (33)
$`\{V_2,V_3,V_4\},\{V_2,V_4,V_5\},\{V_3,V_4,V_5\}\},`$
$`𝒜_0^+`$ $`=`$ $`\{\{V_1,V_2,V_3\},\{V_1,V_2,V_5\},\{V_1,V_3,V_5\},\{V_2,V_3,V_5\},`$ (34)
$`\{V_1,V_4\},\{V_2,V_4\},\{V_3,V_4\},\{V_4,V_5\}\}.`$
Then, the cumulative map $`\psi _{\mathrm{\Gamma }_2}`$ is constructed as follows:
$`V_1`$ $`=`$ $`\psi _{\mathrm{\Gamma }_2}(1)=\{W_4^{(8)},W_6^{(8)},W_7^{(8)},W_8^{(8)}\},`$ (35)
$`V_2`$ $`=`$ $`\psi _{\mathrm{\Gamma }_2}(2)=\{W_3^{(8)},W_5^{(8)},W_7^{(8)},W_8^{(8)}\},`$ (36)
$`V_3`$ $`=`$ $`\psi _{\mathrm{\Gamma }_2}(3)=\{W_2^{(8)},W_5^{(8)},W_6^{(8)},W_8^{(8)}\},`$ (37)
$`V_4`$ $`=`$ $`\psi _{\mathrm{\Gamma }_2}(4)=\{W_1^{(8)},W_2^{(8)},W_3^{(8)},W_4^{(8)}\},`$ (38)
$`V_5`$ $`=`$ $`\psi _{\mathrm{\Gamma }_2}(5)=\{W_1^{(8)},W_5^{(8)},W_6^{(8)},W_7^{(8)}\},`$ (39)
which attains that $`\stackrel{~}{\rho }=\rho ^{}=4`$. On the other hand, the modified cumulative map $`\psi _{\mathrm{\Gamma }_2}^{}`$ is given by
$`V_1`$ $`=`$ $`\psi _{\mathrm{\Gamma }_2}^{}(1)=\{W_1^{(7)},W_9^{(7)}\},`$ (40)
$`V_2`$ $`=`$ $`\psi _{\mathrm{\Gamma }_2}^{}(2)=\{W_2^{(7)},W_8^{(7)}\},`$ (41)
$`V_3`$ $`=`$ $`\psi _{\mathrm{\Gamma }_2}^{}(3)=\{W_3^{(7)},W_7^{(7)}\},`$ (42)
$`V_4`$ $`=`$ $`\psi _{\mathrm{\Gamma }_2}^{}(4)=\{W_4^{(7)},W_6^{(7)},W_7^{(7)},W_8^{(7)},W_9^{(7)}\},`$ (43)
$`V_5`$ $`=`$ $`\psi _{\mathrm{\Gamma }_2}^{}(5)=\{W_5^{(7)},W_6^{(7)}\}.`$ (44)
Observe that the rates of $`\psi _{\mathrm{\Gamma }_2}^{}`$ are given by $`\stackrel{~}{\rho }=13/5,\rho ^{}=5`$. Hence, $`\psi _{\mathrm{\Gamma }_2}^{}`$ gives smaller $`\stackrel{~}{\rho }`$ but larger $`\rho ^{}`$ than $`\psi _{\mathrm{\Gamma }_2}`$. $`\mathrm{}`$
As shown in Examples 5 and 6, the modified cumulative map cannot always overcome the defects of the original cumulative maps. Hence, in the next section, we propose a construction method of multiple assignment maps that can attain the optimal average or worst case coding rates based on integer programming.
## 3 Optimal Multiple Assignment Maps
For a multiple assignment map $`\phi _\mathrm{\Gamma }:\{1,2,\mathrm{},n\}2^{𝑾_{(t,m)}}`$, a set $`𝑨𝑽`$, and $`p\{0,1,\mathrm{},2^n1\}`$, let $`𝑿_p`$ be the subset of $`𝑾_{(t,m)}`$ defined by
$`𝑿_p=\left[{\displaystyle \underset{i:b(p)_i=1}{}}\phi _\mathrm{\Gamma }(i)\right]\left[{\displaystyle \underset{i:b(p)_i=0}{}}\overline{\phi _\mathrm{\Gamma }(i)}\right],`$ (45)
where $`b(p)_i`$ is the $`i`$-th least significant bit in the $`n`$-bit binary representation of $`p`$. For example, in the case of $`p=5`$ and $`n=4`$, it holds that $`b(5)_1=b(5)_3=1`$, and (45) becomes $`𝑿_5=\overline{\phi _\mathrm{\Gamma }(4)}\phi _\mathrm{\Gamma }(3)\overline{\phi _\mathrm{\Gamma }(2)}\phi _\mathrm{\Gamma }(1)`$. Figure 1 is the Venn diagram which shows the relation between $`𝑿_p`$’s and $`\phi _\mathrm{\Gamma }(i)`$’s in the case of $`n=3`$. Since $`\phi _\mathrm{\Gamma }`$ must satisfy (13), it must hold that $`_{i=1}^n\overline{\phi _\mathrm{\Gamma }(i)}=\mathrm{}`$, which implies that $`𝑿_0=\mathrm{}`$. Hence, we consider only $`𝑿_p`$ for $`p=1,2,\mathrm{},2^n1`$ in the following.
Then, it is easy to check that $`𝑿_p`$’s satisfy the following equations for an arbitrary $`n`$ and $`N\stackrel{\text{def}}{=}2^n1`$.
$`𝑿_p𝑿_p^{}`$ $`=`$ $`\mathrm{}\text{if}pp^{}`$ (46)
$`\phi _\mathrm{\Gamma }(i)`$ $`=`$ $`{\displaystyle \underset{p:b(p)_i=1}{}}𝑿_p`$ (47)
$`\mathrm{\Phi }_\mathrm{\Gamma }(𝑨)`$ $`=`$ $`{\displaystyle \underset{V_i𝑨}{}}\phi _\mathrm{\Gamma }(i)={\displaystyle \underset{\genfrac{}{}{0pt}{}{p:b(p)_i=1}{\text{for some}V_i𝑨}}{}}𝑿_p`$ (48)
Letting $`x_p=|𝑿_p|`$, the cardinality of $`\mathrm{\Phi }_\mathrm{\Gamma }(𝑨)`$ is given by
$`|\mathrm{\Phi }_\mathrm{\Gamma }(𝑨)|={\displaystyle \underset{\genfrac{}{}{0pt}{}{p:b(p)_i=1}{\text{for some }V_i𝑨}}{}}x_p,`$ (49)
from (46) and (48).
Now, we describe how to design the optimal multiple assignment map $`\stackrel{~}{\phi }_\mathrm{\Gamma }`$ which attains the minimum average coding rate. Note that, in order to design the multiple assignment map $`\phi _\mathrm{\Gamma }`$ for the set of primitive shares $`𝑾_{(t,m)}`$, we have to determine only $`x_p`$, $`p=1,2,\mathrm{},N`$, and $`t`$, since $`m`$ can be calculated as $`m=_{p=1}^Nx_p`$ from (13) and (49).
Let $`𝒚\stackrel{\text{def}}{=}[t,x_1,x_2,\mathrm{},x_N]`$ be the $`(N+1)`$-dimensional parameter vector to minimize the average coding rate. Furthermore, for an integer $`\mathrm{}`$ and a share set $`𝑨`$, define an $`(N+1)`$-dimensional row vector $`𝒂(\mathrm{};𝑨)\stackrel{\text{def}}{=}[\mathrm{},1(𝑨)_1,1(𝑨)_2,\mathrm{},1(𝑨)_N]`$ where
$`1(𝑨)_p=\{\begin{array}{cc}1\hfill & \text{if}b(p)_i=1\text{for some}V_i𝑨\hfill \\ 0\hfill & \text{otherwise.}\hfill \end{array}`$ (52)
Then, since (49) can be represented by inner product as $`|\mathrm{\Phi }_\mathrm{\Gamma }(𝑨)|=𝒂(0;𝑨)𝒚^T`$ where superscript $`T`$ means the transpose of vector $`𝒚`$, the inequalities in the constraints (11) and (12) can be represented by $`𝒂(0;𝑨)𝒚^Tt`$, and $`𝒂(0;𝑨)𝒚^Tt1`$, respectively. Therefore, these constraints can be expressed as
$`𝒂(1;𝑨)𝒚^T`$ $``$ $`0\text{if}𝑨𝒜_1^{},`$ (53)
$`𝒂(1;𝑨)𝒚^T1`$ $``$ $`0\text{if}𝑨𝒜_0^+,`$ (54)
respectively. Furthermore, denoting the Hamming weight in the binary representation of $`p`$ by $`h_p`$, it holds from (47) that
$`{\displaystyle \underset{i=1}{\overset{n}{}}}|\phi _\mathrm{\Gamma }(i)|={\displaystyle \underset{i=1}{\overset{n}{}}}{\displaystyle \underset{p:b(p)_i=1}{}}x_p={\displaystyle \underset{p=1}{\overset{N}{}}}h_px_p=𝒉𝒚^T,`$ (55)
where $`𝒉=[h_0,h_1,\mathrm{},h_N]^{N+1}`$. Hence, the average coding rate $`\stackrel{~}{\rho }`$ in (14) is given by $`(1/n)𝒉𝒚^T`$ which we want to minimize.
We note here that $`𝒂(;)`$ and $`𝒉`$ do not depend on the multiple assignment map $`\phi _\mathrm{\Gamma }`$, and hence, summarizing (52)–(55), we can formulate the integer programming problem IP$`{}_{\stackrel{~}{\rho }}{}^{}(\mathrm{\Gamma })`$ that minimizes the average coding rate $`\stackrel{~}{\rho }`$ under the constraints of (11) and (12) as follows:
| IP$`{}_{\stackrel{~}{\rho }}{}^{}(\mathrm{\Gamma })`$ | | | | | |
| --- | --- | --- | --- | --- | --- |
| minimize | $`𝒉𝒚^T`$ | | | | |
| subject to | $`𝒂(1;𝑨)𝒚^T`$ | $``$ | $`0`$ | for | $`𝑨𝒜_1^{}`$ |
| | $`𝒂(1;𝑨)𝒚^T`$ | $``$ | $`1`$ | for | $`𝑨𝒜_0^+`$ |
| | $`𝒚`$ | $``$ | $`\mathrm{𝟎}`$ | | |
The optimal multiple assignment map $`\stackrel{~}{\phi }_\mathrm{\Gamma }`$ that attains the minimum average coding rate can be constructed as follows. First, let $`\stackrel{~}{𝒚}=[\stackrel{~}{t},\stackrel{~}{x}_1,\stackrel{~}{x}_2,\mathrm{},\stackrel{~}{x}_N]`$ be the minimizers of the integer programming problem IP$`{}_{\stackrel{~}{\rho }}{}^{}(\mathrm{\Gamma })`$, and we use the $`(\stackrel{~}{t},\stackrel{~}{m})`$-threshold SSS with primitive shares $`𝑾_{(\stackrel{~}{t},\stackrel{~}{m})}=\{W_1^{(\stackrel{~}{t})},W_2^{(\stackrel{~}{t})},\mathrm{},W_{\stackrel{~}{m}}^{(\stackrel{~}{t})}\}`$ for secret $`S`$ where $`\stackrel{~}{m}`$ can be calculated from $`\stackrel{~}{m}=_{p=1}^N\stackrel{~}{x}_p`$. Then, for each $`p`$, we can assign $`\stackrel{~}{x}_p`$ different primitive shares of $`𝑾_{(\stackrel{~}{t},\stackrel{~}{m})}`$ to $`𝑿_p`$ that satisfies $`|𝑿_p|=\stackrel{~}{x}_p`$ and (46). Finally, the multiple assignment map $`\stackrel{~}{\phi }_\mathrm{\Gamma }`$ is obtained by (47).
Next, we consider the integer programming problem IP$`{}_{\rho ^{}}{}^{}(\mathrm{\Gamma })`$ that minimizes the worst coding rate $`\rho ^{}`$. Let $`M`$ be the maximal number of assigned primitive shares among all $`V_i`$, $`i=1,2,\mathrm{},n`$. Then, it holds that $`|\phi _\mathrm{\Gamma }(i)|M`$ for all $`i=1,2,\mathrm{},n`$, and the minimization of $`M`$ attains the optimal worst coding rate. Now, let $`𝒛`$ be the $`(N+2)`$-dimensional parameter vector defined by $`𝒛\stackrel{\text{def}}{=}[M,t,x_1,x_2,\mathrm{},x_N]`$. Then, it holds that $`M=𝒆𝒛^T`$ where $`𝒆`$ is the $`(N+2)`$-dimensional row vector defined by $`𝒆\stackrel{\text{def}}{=}[1,0,0,\mathrm{},0]`$. Furthermore, by defining $`𝒃(\mathrm{},\mathrm{}^{};𝑨)\stackrel{\text{def}}{=}[\mathrm{},\mathrm{}^{},1(𝑨)_1,1(𝑨)_2,\mathrm{},1(𝑨)_N]`$ where $`1(𝑨)_p`$ is defined by (52), the number of primitive shares assigned to a share set $`𝑨𝑽`$ can be expressed as $`𝒃(0,0;𝑨)𝒛^T`$. Hence, in the same way as IP$`{}_{\stackrel{~}{\rho }}{}^{}(\mathrm{\Gamma })`$, the integer programming problem IP$`{}_{\rho ^{}}{}^{}(\mathrm{\Gamma })`$ that minimizes the worst coding rate $`\rho ^{}`$ can be formulated as follows:
| IP$`{}_{\rho ^{}}{}^{}(\mathrm{\Gamma })`$ | | | | | |
| --- | --- | --- | --- | --- | --- |
| minimize | $`𝒆𝒛^T`$ | | | | |
| subject to | $`𝒃(0,1;𝑨)𝒛^T`$ | $``$ | $`0`$ | for | $`𝑨𝒜_1^{}`$ |
| | $`𝒃(0,1;𝑨)𝒛^T`$ | $``$ | $`1`$ | for | $`𝑨𝒜_0^+`$ |
| | $`𝒃(1,0;\{V\})𝒛^T`$ | $``$ | $`0`$ | for | $`V𝑽`$ |
| | $`𝒛`$ | $``$ | $`\mathrm{𝟎}`$ | | |
The multiple assignment map $`\phi _\mathrm{\Gamma }^{}`$ attaining the minimum $`\rho ^{}`$ can also be constructed from the obtained minimizer in the same way as the construction of $`\stackrel{~}{\phi }_\mathrm{\Gamma }`$.
###### Remark 7
Actually, in SSSs, we can assume without loss of generality that $`x_N=0`$, i.e., $`𝑿_N=_{i=1}^n\phi _\mathrm{\Gamma }(i)=\mathrm{}`$ because it is not necessary to consider the set of primitive shares commonly contained in every share. Hence, the vectors in integer programming problems IP$`{}_{\stackrel{~}{\rho }}{}^{}(\mathrm{\Gamma })`$ and IP$`{}_{\rho ^{}}{}^{}(\mathrm{\Gamma })`$ can be reduced to $`N`$-dimensional and $`(N+1)`$-dimensional vectors, respectively. However, $`x_N=0`$ does not hold generally in the case of ramp SS schemes, which is described in Remark 20 in Section 5.2. $`\mathrm{}`$
###### Example 8
For the access structure $`\mathrm{\Gamma }_1`$ defined by (18) and (19) in Example 2, the integer programming problem IP$`{}_{\stackrel{~}{\rho }}{}^{}(\mathrm{\Gamma }_1)`$ can be formulated as follows:
| IP$`{}_{\stackrel{~}{\rho }}{}^{}(\mathrm{\Gamma }_1)`$ | |
| --- | --- |
| minimize | $`x_1+x_2+2x_3+x_4+2x_5+2x_6+3x_7+x_8+2x_9+2x_{10}`$ |
| | $`+3x_{11}+2x_{12}+3x_{13}+3x_{14}`$ |
| subject to | $`t+x_1+x_2+x_3+x_4+x_5+x_6+x_7+x_9`$ |
| | $`+x_{10}+x_{11}+x_{12}+x_{13}+x_{14}`$ $``$ $`0`$ $`t+x_1+x_3+x_5+x_7+x_8+x_9+x_{10}+x_{11}+x_{12}+x_{13}+x_{14}`$ $``$ $`0`$ $`t+x_2+x_3+x_6+x_7+x_8+x_9+x_{10}+x_{11}+x_{12}+x_{13}+x_{14}`$ $``$ $`0`$ $`t+x_4+x_5+x_6+x_7+x_8+x_9+x_{10}+x_{11}+x_{12}+x_{13}+x_{14}`$ $``$ $`0`$ $`tx_1x_2x_3x_5x_6x_7x_9x_{10}x_{11}x_{13}x_{14}`$ $``$ $`1`$ $`tx_1x_3x_4x_5x_6x_7x_9x_{11}x_{12}x_{13}x_{14}`$ $``$ $`1`$ $`tx_2x_3x_4x_5x_6x_7x_{10}x_{11}x_{12}x_{13}x_{14}`$ $``$ $`1`$ $`tx_8x_9x_{10}x_{11}x_{12}x_{13}x_{14}`$ $``$ $`1`$ $`x_p`$ $``$ $`0,p=1,2,\mathrm{},14`$ |
By solving the above IP$`{}_{\stackrel{~}{\rho }}{}^{}(\mathrm{\Gamma }_1)`$, we obtain that the value of the objective function is $`5`$, which is attained by the following minimizers:
$`\stackrel{~}{t}=3,\stackrel{~}{x}_1`$ $`=`$ $`\stackrel{~}{x}_2=\stackrel{~}{x}_4=1,\stackrel{~}{x}_8=2,\stackrel{~}{x}_i=0\mathrm{for}i=3,5,6,7,9,10,\mathrm{},14,`$ (56)
Hence, $`\stackrel{~}{m}`$ is given by $`\stackrel{~}{m}=_{p=1}^{14}\stackrel{~}{x}_p=5`$, and $`𝑿_p`$’s become
$`𝑿_1`$ $`=`$ $`\left\{W_1^{(3)}\right\},𝑿_2=\left\{W_2^{(3)}\right\},𝑿_4=\left\{W_3^{(3)}\right\},𝑿_8=\{W_4^{(3)},W_5^{(3)}\},`$ (57)
where $`𝑾_{(3,5)}=\{W_1^{(3)},W_2^{(3)},\mathrm{},W_5^{(3)}\}`$. Finally, from (47), $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_1}`$ is constructed as
$`V_1`$ $`=`$ $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_1}(1)=\left\{W_1^{(3)}\right\},`$ (58)
$`V_2`$ $`=`$ $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_1}(2)=\left\{W_2^{(3)}\right\},`$ (59)
$`V_3`$ $`=`$ $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_1}(3)=\left\{W_3^{(3)}\right\},`$ (60)
$`V_4`$ $`=`$ $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_1}(4)=\{W_4^{(3)},W_5^{(3)}\}.`$ (61)
In this case, we have that $`\stackrel{~}{\rho }=5/4`$ and $`\rho ^{}=2`$. The integer programming problem $`\mathrm{IP}_\rho ^{}(\mathrm{\Gamma }_1)`$ derives the same solutions as (56), and hence, it holds that $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_1}=\phi _{\mathrm{\Gamma }_1}^{}`$ in this example. Recall that the cumulative map $`\psi _{\mathrm{\Gamma }_1}`$ attains the coding rates $`\stackrel{~}{\rho }=9/4`$ and $`\rho ^{}=3`$, and the modified cumulative map $`\psi _{\mathrm{\Gamma }_1}^{}`$ attains $`\stackrel{~}{\rho }=5/2`$ and $`\rho ^{}=4`$. Hence, $`\phi _{\mathrm{\Gamma }_1}`$ can attain smaller coding rates compared with $`\psi _{\mathrm{\Gamma }_1}`$ and $`\psi _{\mathrm{\Gamma }_1}^{}`$. $`\mathrm{}`$
###### Example 9
For the access structure $`\mathrm{\Gamma }_2`$ defined by (33) and (34) in Example 6, we can obtain the following multiple assignment map by solving the integer programming problem $`\mathrm{IP}_{\stackrel{~}{\rho }}(\mathrm{\Gamma }_2)`$.
$`V_1`$ $`=`$ $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_2}(1)=\left\{W_1^{(4)}\right\},`$ (62)
$`V_2`$ $`=`$ $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_2}(2)=\left\{W_2^{(4)}\right\},`$ (63)
$`V_3`$ $`=`$ $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_2}(3)=\left\{W_3^{(4)}\right\},`$ (64)
$`V_4`$ $`=`$ $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_2}(4)=\{W_4^{(4)},W_5^{(4)}\},`$ (65)
$`V_5`$ $`=`$ $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_2}(5)=\left\{W_6^{(4)}\right\},`$ (66)
where $`W_i^{(4)}𝑾_{(4,6)}`$. Then, it holds that $`\stackrel{~}{\rho }=6/5`$ and $`\rho ^{}=2`$. Furthermore, it holds that $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_2}=\phi _{\mathrm{\Gamma }_2}^{}`$ in this access structure. Recall again that the cumulative map $`\psi _{\mathrm{\Gamma }_2}`$ attains the coding rates $`\stackrel{~}{\rho }=\rho ^{}=4`$, and the modified cumulative map $`\psi _{\mathrm{\Gamma }_2}^{}`$ attains $`\stackrel{~}{\rho }=13/5`$ and $`\rho ^{}=5`$. Hence, $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_2}`$ is more efficient than $`\psi _{\mathrm{\Gamma }_2}`$ and $`\psi _{\mathrm{\Gamma }_2}^{}`$. $`\mathrm{}`$
Since any access structure can be realized by the cumulative map (and the modified cumulative map), there exists at least one multiple assignment map for any access structure. Therefore, the next theorem holds obviously.
###### Theorem 10
For any access structure $`\mathrm{\Gamma }`$ that satisfies monotonicity (3) and (4), the integer programming problems IP$`{}_{\stackrel{~}{\rho }}{}^{}(\mathrm{\Gamma })`$ and IP$`{}_{\rho ^{}}{}^{}(\mathrm{\Gamma })`$ always have at least one feasible solution, and hence, there exists the optimal multiple assignment map. $`\mathrm{}`$
We note that the integer programming problems are NP-hard, and hence, the proposed algorithms may take much time in solving for large $`n`$ $`(=|𝑽|)`$. But, in the case that $`n`$ is not large, the solution is obtained quickly. For instance, in the case of IP$`{}_{\rho }{}^{}(\mathrm{\Gamma }_3)`$ in Example 11 with $`n=6`$, it can be solved within $`0.1`$ seconds by a notebook computer.
###### Example 11
Consider the following access structure $`\mathrm{\Gamma }_3`$:
$`𝒜_1^{}`$ $`=`$ $`\{\{V_1,V_3,V_4,V_5\},\{V_1,V_3,V_5,V_6\},\{V_1,V_4,V_5,V_6\},\{V_3,V_4,V_5,V_6\},\{V_1,V_2,V_3\},\{V_1,V_2,V_5\},`$
$`\{V_1,V_2,V_6\},\{V_2,V_3,V_4\},\{V_2,V_3,V_5\},\{V_2,V_3,V_6\},\{V_2,V_4,V_5\},\{V_2,V_4,V_6\},\{V_2,V_5,V_6\}\},`$
$`𝒜_0^+`$ $`=`$ $`\{\{V_1,V_3,V_4,V_6\},\{V_1,V_2,V_4\},\{V_1,V_3,V_5\},\{V_1,V_4,V_5\},\{V_1,V_5,V_6\},\{V_3,V_4,V_5\},`$ (68)
$`\{V_3,V_5,V_6\},\{V_4,V_5,V_6\},\{V_2,V_3\},\{V_2,V_5\},\{V_2,V_6\}\}.`$
Then, we obtain the following multiple assignment map by solving IP$`{}_{\stackrel{~}{\rho }}{}^{}\left(\mathrm{\Gamma }_3\right)`$.
$`V_1`$ $`=`$ $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_3}(1)=\{W_1^{(6)},W_2^{(6)}\},`$ (69)
$`V_2`$ $`=`$ $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_3}(2)=\{W_1^{(6)},W_3^{(6)},W_4^{(6)},W_5^{(6)}\},`$ (70)
$`V_3`$ $`=`$ $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_3}(3)=\left\{W_6^{(6)}\right\},`$ (71)
$`V_4`$ $`=`$ $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_3}(4)=\{W_2^{(6)},W_5^{(6)}\},`$ (72)
$`V_5`$ $`=`$ $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_3}(5)=\{W_3^{(6)},W_7^{(6)}\},`$ (73)
$`V_6`$ $`=`$ $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_3}(6)=\left\{W_8^{(6)}\right\},`$ (74)
where $`W_i^{(6)}𝑾_{(6,8)}`$. $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_3}`$ attains that $`\stackrel{~}{\rho }=2`$ and $`\rho ^{}=4`$. On the other hand, the cumulative map for the access structure $`\mathrm{\Gamma }_3`$ are given by
$`V_1`$ $`=`$ $`\psi _{\mathrm{\Gamma }_3}(1)=\{W_6^{(11)},W_7^{(11)},W_8^{(11)},W_9^{(11)},W_{10}^{(11)},W_{11}^{(11)}\},`$ (75)
$`V_2`$ $`=`$ $`\psi _{\mathrm{\Gamma }_3}(2)=\{W_1^{(11)},W_3^{(11)},W_4^{(11)},W_5^{(11)},W_6^{(11)},W_7^{(11)},W_8^{(11)}\},`$ (76)
$`V_3`$ $`=`$ $`\psi _{\mathrm{\Gamma }_3}(3)=\{W_2^{(11)},W_4^{(11)},W_5^{(11)},W_8^{(11)},W_{10}^{(11)},W_{11}^{(11)}\},`$ (77)
$`V_4`$ $`=`$ $`\psi _{\mathrm{\Gamma }_3}(4)=\{W_3^{(11)},W_5^{(11)},W_7^{(11)},W_9^{(11)},W_{10}^{(11)},W_{11}^{(11)}\},`$ (78)
$`V_5`$ $`=`$ $`\psi _{\mathrm{\Gamma }_3}(5)=\{W_1^{(11)},W_2^{(11)},W_9^{(11)},W_{11}^{(11)}\},`$ (79)
$`V_6`$ $`=`$ $`\psi _{\mathrm{\Gamma }_3}(6)=\{W_2^{(11)},W_3^{(11)},W_4^{(11)},W_6^{(11)},W_9^{(11)},W_{10}^{(11)}\},`$ (80)
where $`W_i^{(11)}𝑾_{(11,11)}`$. $`\psi _{\mathrm{\Gamma }_3}`$ has $`\stackrel{~}{\rho }=35/6`$ and $`\rho ^{}=7`$. Furthermore, the modified cumulative map for $`\mathrm{\Gamma }_3`$ requires $`(12,15)`$-threshold SSS and has $`\stackrel{~}{\rho }=5`$ and $`\rho ^{}=9`$. $`\mathrm{}`$
Next, we clarify what kind of access structure can be realized as an ideal SSS by the multiple assignment map.
###### Theorem 12
For an access structure $`\mathrm{\Gamma }`$, the SSS constructed by the optimal multiple assignment map is ideal, i.e., $`\rho _i=1`$ for all $`i`$, if and only if $`𝒜_1^{}`$ of $`\mathrm{\Gamma }`$ can be represented by
$`𝒜_1^{}`$ $`=`$ $`{\displaystyle \underset{\genfrac{}{}{0pt}{}{\{j_1,j_2,\mathrm{},j_t\}}{\{1,2,\mathrm{},m\}}}{}}\left\{𝑨_{j_1}\times 𝑨_{j_2}\times \mathrm{}\times 𝑨_{j_t}\right\},`$ (81)
where $`t`$ is a positive integer and $`\{𝑨_1,𝑨_2,\mathrm{},𝑨_m\}`$ is a partition of $`𝑽`$ which satisfies
$`{\displaystyle \underset{j=1}{\overset{m}{}}}𝑨_j`$ $`=`$ $`𝑽,`$ (82)
$`𝑨_j`$ $``$ $`\mathrm{}\text{for}j=1,2,\mathrm{},m,`$ (83)
$`𝑨_j𝑨_j^{}`$ $`=`$ $`\mathrm{}\text{if}jj^{}.`$ (84)
$`\mathrm{}`$
Proof of Theorem 12: If there exists a partition $`\{𝑨_1,𝑨_2,\mathrm{},𝑨_m\}`$ satisfying (81)–(84) for the access structure $`\mathrm{\Gamma }`$, the ideal SSS can be obtained by letting
$`\phi _\mathrm{\Gamma }(i)=W_j^{(t)}\text{if}V_i𝑨_j`$ (85)
for each $`i=1,2,\mathrm{},n`$. Next, we show the necessity of (81)–(84). Suppose that a certain $`\phi _\mathrm{\Gamma }(i)`$ attains $`\rho _i=1`$ for all $`i`$. Then, define each $`𝑨_j`$ as
$`𝑨_j\stackrel{\text{def}}{=}\mathrm{\Phi }_\mathrm{\Gamma }^1\left(\left\{W_j^{(t)}\right\}\right),j=1,2,\mathrm{},m,`$ (86)
for $`j=1,2,\mathrm{},m`$ where $`\mathrm{\Phi }_\mathrm{\Gamma }^1:2^{𝑾_{(t,m)}}2^𝑽`$ is the inverse map of $`\mathrm{\Phi }_\mathrm{\Gamma }(𝑨)\stackrel{\text{def}}{=}_{i:V_i𝑨}\phi _\mathrm{\Gamma }(i)`$. Then, it is easy to see that $`𝑨_j`$’s satisfy (81), (82) and (83). Next, we prove that $`𝑨_j`$’s defined by (86) satisfy (84). Assume that there exist $`𝑨_j`$ and $`𝑨_j^{}`$, $`jj^{}`$, not satisfying (84). Then, there exists a share $`V_i𝑨_j𝑨_j^{}`$. This means that $`\phi _\mathrm{\Gamma }(i)\{W_j^{(t)},W_j^{}^{(t)}\}`$, which contradicts $`\rho _i=|\phi _\mathrm{\Gamma }(i)|=1`$. Hence, $`\{𝑨_1,𝑨_2,\mathrm{},𝑨_m\}`$ must be a partition of $`𝑽`$ satisfying (81)–(84). $`\mathrm{}`$
In the case of $`t=2`$, it is known that an access structure $`\mathrm{\Gamma }`$ can be realized by an ideal SSS if and only if $`\mathrm{\Gamma }`$ can be represented by a complete multipartite graph . We note that this condition coincides with (81)–(84) in this case. Furthermore, in the case that $`|𝑨_j|=1`$ for $`j=1,2,\mathrm{},m`$, the access structure coincides with the $`(t,m)`$-threshold access structure. Hence, if $`\mathrm{\Gamma }`$ is the $`(k,n)`$-threshold access structure, the multiple assignment maps obtained from the integer programming problems IP$`{}_{\stackrel{~}{\rho }}{}^{}(\mathrm{\Gamma })`$ and IP$`{}_{\rho ^{}}{}^{}(\mathrm{\Gamma })`$ obviously satisfy that $`|\stackrel{~}{\phi }_\mathrm{\Gamma }(i)|=|\phi _\mathrm{\Gamma }^{}(i)|=1`$ for all $`i`$.
We note that any access structures not satisfying (81)–(84) must have $`\stackrel{~}{\rho }>1`$ and $`\rho ^{}2`$ if the multiple assignment map is used. But, an access structure not satisfying (81)–(84) might be realized as an ideal SSS if we use another construction method. For example, refer .
In this paper, we assume that every share is significant. But, if there exist vacuous shares in the access structure $`\mathrm{\Gamma }`$, it is cumbersome to check whether each share is significant or vacuous. From Remark 1, the optimal multiple assignment map $`\stackrel{~}{\phi }_\mathrm{\Gamma }`$ attaining the minimum average coding rate must satisfy that $`|\stackrel{~}{\phi }_\mathrm{\Gamma }(i)|=0`$ for any vacuous share $`V_i`$. On the other hand, it clearly holds that $`|\phi _\mathrm{\Gamma }(i)|1`$ for every significant share $`V_i`$ since $`\rho _i1`$ holds for any significant share. Hence, by solving the integer programming problem IP$`{}_{\stackrel{~}{\rho }}{}^{}(\mathrm{\Gamma })`$, we can also know whether a share is significant or vacuous.
## 4 Multiple Assignment Maps for Incomplete Access Structures
In the previous sections, we considered how to construct a SSS for a complete general access structure $`\mathrm{\Gamma }=\{𝒜_1,𝒜_0\}`$. But in practice, it may be cumbersome to specify whether each subset of $`𝑽`$ is a qualified set or a forbidden set because the number of the subsets is $`2^n`$. Hence, a method is proposed in to construct a SSS for the case such that some subsets of $`𝑽`$ are not specified as qualified nor forbidden sets.
###### Theorem 13 ()
Let $`\mathrm{\Gamma }^{\mathrm{}}=\{𝒜_1^{\mathrm{}},𝒜_0^{\mathrm{}}\}`$ be an incomplete access structure, which has $`𝒜_1^{\mathrm{}}𝒜_1^{\mathrm{}}2^𝑽`$. Then, there exists a complete access structure $`\mathrm{\Gamma }=\{𝒜_1,𝒜_0\}`$ such that
$`𝒜_1^{\mathrm{}}`$ $``$ $`𝒜_1,`$ (87)
$`𝒜_0^{\mathrm{}}`$ $``$ $`𝒜_0,`$ (88)
if and only if it holds that for any $`𝑨𝒜_1^{\mathrm{}}`$ and $`𝑩𝒜_0^{\mathrm{}}`$,
$`𝑨𝑩.`$ (89)
$`\mathrm{}`$
In case that (89) is satisfied, the SSS satisfying the incomplete access structure $`\mathrm{\Gamma }^{\mathrm{}}=\{𝒜_1^{\mathrm{}},𝒜_0^{\mathrm{}}\}`$ can be realized by applying the cumulative map to the complete access structure $`\mathrm{\Gamma }=\{𝒜_1,𝒜_0\}`$. In fact, for the access structure $`\mathrm{\Gamma }^{\mathrm{}}=\{𝒜_1^{\mathrm{}},𝒜_0^{\mathrm{}}\}`$, a SSS is constructed in by a cumulative map $`\psi _\mathrm{\Gamma }^{\mathrm{}}(i)=_{j:V_i𝑭_j}\{W_j^{(t)}\}`$ for $`𝒜_0^\mathrm{}+=\{𝑭_1,𝑭_2,\mathrm{},𝑭_m\}`$. This construction corresponds to the case that
$`𝒜_0^+=𝒜_0^\mathrm{}+\text{and}𝒜_1=2^𝑽𝒜_0.`$ (90)
However, $`\psi _\mathrm{\Gamma }^{\mathrm{}}`$ is not efficient generally because $`\psi _\mathrm{\Gamma }^{\mathrm{}}`$ is a cumulative map, which is inefficient as described in Section 2.2. Furthermore, even if the cumulative map can attain the optimal coding rates for the access structure given by (90), the access structure may not be optimal among all the complete access structures $`\mathrm{\Gamma }=\{𝒜_1,𝒜_0\}`$ satisfying (87) and (88) for given $`\mathrm{\Gamma }^{\mathrm{}}=\{𝒜_1^{\mathrm{}},𝒜_0^{\mathrm{}}\}`$.
In our construction based on integer programming, the optimal multiple assignment map for the incomplete access structure $`\mathrm{\Gamma }^{\mathrm{}}=\{𝒜_1^{\mathrm{}},𝒜_0^\mathrm{}+\}`$ can easily be obtained by applying IP$`{}_{\stackrel{~}{\rho }}{}^{}(\mathrm{\Gamma })`$ or IP$`{}_{\rho ^{}}{}^{}(\mathrm{\Gamma })`$ directly to $`\mathrm{\Gamma }^{\mathrm{}}`$.
###### Example 14
Let us consider the following access structure $`\mathrm{\Gamma }_3^{\mathrm{}}=\{𝒜_1^{\mathrm{}},𝒜_0^{\mathrm{}}\}`$:
$`𝒜_1^{\mathrm{}}`$ $`=`$ $`\{\{V_1,V_4,V_5,V_6\},\{V_1,V_2,V_5\},\{V_1,V_2,V_6\},\{V_2,V_3,V_6\},\{V_2,V_4,V_6\}\},`$ (91)
$`𝒜_0^{\mathrm{}}`$ $`=`$ $`\{\{V_1,V_3,V_4,V_6\},\{V_1,V_3,V_5\},\{V_1,V_5,V_6\},\{V_3,V_4,V_5\},\{V_4,V_5,V_6\},\{V_2,V_5\}\},`$ (92)
Note that $`𝒜_1^{\mathrm{}}`$ and $`𝒜_0^{\mathrm{}}`$ satisfy $`𝒜_1^{\mathrm{}}𝒜_1^{}`$ and $`𝒜_0^{\mathrm{}}𝒜_0^+`$ for $`\mathrm{\Gamma }_3=\{𝒜_1,𝒜_0\}`$, which is defined by (LABEL:tochi1.eq) and (68) in Example 11. Then, by solving $`\mathrm{IP}_{\stackrel{~}{\rho }}(\mathrm{\Gamma }_3^{\mathrm{}})`$, we obtain the following multiple assignment map.
$`V_1`$ $`=`$ $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_3^{\mathrm{}}}(1)=\left\{W_1^{(4)}\right\},`$ (93)
$`V_2`$ $`=`$ $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_3^{\mathrm{}}}(2)=\{W_2^{(4)},W_3^{(4)}\},`$ (94)
$`V_3`$ $`=`$ $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_3^{\mathrm{}}}(3)=\left\{W_4^{(4)}\right\},`$ (95)
$`V_4`$ $`=`$ $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_3^{\mathrm{}}}(4)=\left\{W_4^{(4)}\right\},`$ (96)
$`V_5`$ $`=`$ $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_3^{\mathrm{}}}(5)=\left\{W_5^{(4)}\right\},`$ (97)
$`V_6`$ $`=`$ $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_3^{\mathrm{}}}(6)=\left\{W_6^{(4)}\right\},`$ (98)
where $`W_i^{(4)}𝑾_{(4,6)}`$, and it holds that $`\stackrel{~}{\rho }=7/6`$ and $`\rho ^{}=2`$. If we apply the cumulative map to $`\mathrm{\Gamma }_3^{\mathrm{}}`$, $`\psi _{\mathrm{\Gamma }_3^{\mathrm{}}}`$ is constructed from the $`(6,6)`$-threshold scheme, and it has $`\stackrel{~}{\rho }=3`$ and $`\rho ^{}=5`$. $`\mathrm{}`$
Similarly to the complete SSS, vacuous shares $`V_i`$ in $`\mathrm{\Gamma }^{\mathrm{}}=\{𝒜_1^{\mathrm{}},𝒜_0^{\mathrm{}}\}`$ can be detected by checking $`|\phi _\mathrm{\Gamma }^{\mathrm{}}(i)|=0`$ for the solution of the IP$`{}_{\stackrel{~}{\rho }}{}^{}(\mathrm{\Gamma }^{\mathrm{}})`$.
## 5 Ramp SSSs with General Access Structures
The coding rate $`\rho _i`$ must satisfy $`\rho _i1`$ for any significant share $`V_i`$ in the case that the access structure consists of $`𝒜_1`$ and $`𝒜_0`$, i.e., every subset $`𝑨𝑽`$ is classified into either qualified sets or forbidden sets. But, in the case of ramp access structures such that some subsets of $`𝑽`$ are allowed to have intermediate properties between the qualified and forbidden sets, it is possible to decrease the coding rate $`\rho _i`$ to less than 1. The SSSs having the ramp access structure are called ramp schemes . In this section, we treat the construction of ramp SSSs based on the multiple assignment maps. We consider only the minimum average coding rate in this section. But, for the minimum worst coding rate, integer programming can be formulated in a similar way.
### 5.1 Preliminaries for Ramp Schemes
First, let us review the definition of ramp SSSs. Suppose that $`L+1`$ families $`𝒜_j2^𝑽`$, $`j=0,1,\mathrm{},L`$, satisfy the following.
$`H(S|𝑨)={\displaystyle \frac{Lj}{L}}H(S),\text{for any }𝑨𝒜_j`$ (99)
Equation (99) implies that the secret $`S`$ leaks out from a set $`𝑨𝒜_j`$ with the amount of $`(j/L)H(S)`$. Especially, $`S`$ can be decrypted completely from any $`𝑨𝒜_L`$, and any $`𝑨𝒜_0`$ leaks out no information of $`S`$. Note that, in the case of $`L=1`$, the ramp SSS reduces to the SSS treated in Sections 2–4, and hence, the ramp SSS can be considered as an extension of the ordinal SSS. To distinguish the ordinal SSSs from ramp SSSs, the ordinal SSSs are called the perfect SSSs. We call $`\mathrm{\Gamma }^R=\{𝒜_0,𝒜_1,\mathrm{},𝒜_L\}`$ the access structure of the ramp SSS with $`L+1`$ levels. Without loss of generality, we can assume that $`_{j=0}^L𝒜_j=2^𝑽`$ and $`𝒜_j𝒜_j^{}=\mathrm{}`$ for $`jj^{}`$, although incomplete access structures with $`_{j=0}^L𝒜_j2^𝑽`$ can be treated in the same way as in Section 4.
For example, the access structure of $`(k,L,n)`$-threshold ramp SSS is defined as follows:
$`𝒜_0`$ $`=`$ $`\{𝑨2^𝑽:0|𝑨|kL\},`$ (100)
$`𝒜_j`$ $`=`$ $`\{𝑨2^𝑽:|𝑨|=kL+j\},\text{for }1jL1,`$ (101)
$`𝒜_L`$ $`=`$ $`\{𝑨2^𝑽:k|𝑨|n\}.`$ (102)
In ramp SSSs, a significant share can also be defined in the same way as the perfect SSSs shown in Section 2.1. A share $`V_i𝑽`$ is called significant if there exists a share set $`𝑨2^𝑽`$ such that $`𝑨\{V_i\}𝒜_j`$ and $`𝑨𝒜_j^{}`$ with $`j>j^{}`$. Then, a non-significant share $`V_i^{}`$ satisfies that $`𝑨\{V_i^{}\}𝒜_j`$ for any share set $`𝑨𝒜_j`$, $`j=0,1,\mathrm{},L`$. Furthermore, if a non-significant share $`V_i^{}`$ satisfies $`\{V_i^{}\}𝒜_0`$, $`V_i^{}`$ plays no roll in the ramp SSS, and hence, we call $`V_i^{}`$ a vacuous share. However, there exists a ramp scheme such that $`𝒜_0=\mathrm{}`$ and a non-significant share satisfy $`\{V_i\}𝒜_j`$ for some $`j1`$. This case implies that $`H(V_i^{})H(S)/L`$, and $`H(V_i^{}|V)=0`$ for any $`V𝑽`$, i.e., a non-significant $`V_i^{}`$ is included in every share. Therefore, we call such a non-significant share $`V_i^{}`$ a common share.
###### Remark 15
It is known that for any access structure with $`L+1`$ levels, the coding rate $`\rho _i`$ must satisfy $`\rho _i1/L`$ for any significant share $`V_i`$ . Especially, in the case of $`(k,L,n)`$-threshold SSSs, the optimal ramp SSS attaining $`\rho _i=1/L`$ for all $`i`$ can easily be constructed . Any common share $`V_i`$ must also satisfy that $`\rho _i1/L`$. On the other hand, in the same way as Remark 1 for the perfect SSSs, each vacuous share $`V_i`$ can be realized as $`\rho _i=0`$ for any access structure. Furthermore, if there exists a vacuous share with $`\rho _i>0`$, the average coding rate can be reduced by setting $`\rho _i=0`$ without changing all the significant and the common shares. $`\mathrm{}`$
Letting $`\stackrel{ˇ}{𝒜}_j\stackrel{\text{def}}{=}_{\mathrm{}=j}^L𝒜_{\mathrm{}}`$ and $`\widehat{𝒜}_j\stackrel{\text{def}}{=}_{\mathrm{}=1}^j𝒜_{\mathrm{}}`$, for $`j=0,1,\mathrm{},L`$, the monotonicity in (3) and (4) are extended as follows:
$`𝑨\stackrel{ˇ}{𝒜}_j𝑨^{}\stackrel{ˇ}{𝒜}_j\text{for all}𝑨^{}𝑨`$ (103)
$`𝑨\widehat{𝒜}_j𝑨^{}\widehat{𝒜}_j\text{for all}𝑨^{}𝑨`$ (104)
Therefore, the minimal and maximal families of the access structure, $`\mathrm{\Gamma }^R=\{𝒜_0^{},𝒜_1^{},\mathrm{},𝒜_L^{}\}`$ and $`\mathrm{\Gamma }^{R+}=\{𝒜_0^+,𝒜_1^+,\mathrm{},𝒜_L^+\}`$, respectively, can be defined as
$`𝒜_j^{}`$ $`=`$ $`\{𝑨𝒜_j:𝑨\{V\}\stackrel{ˇ}{𝒜}_j\text{for any}V𝑨\},`$ (105)
$`𝒜_j^+`$ $`=`$ $`\{𝑨𝒜_j:𝑨\{V\}\widehat{𝒜}_j\text{for any}V2^𝑽𝑨\}.`$ (106)
Then, the following theorem holds.
###### Theorem 16 ()
A ramp SSS with access structure $`\mathrm{\Gamma }^R=\{𝒜_0,𝒜_1,\mathrm{},𝒜_L\}`$ can be constructed if and only if $`\stackrel{ˇ}{𝒜}_j`$ (or $`\widehat{𝒜}_j`$) satisfies the monotonicity (103) (or (104)) for all $`j=1,2,\mathrm{},L`$. $`\mathrm{}`$
In Theorem 16, the necessity of the condition is obvious, and the sufficiency is established by the next construction.
###### Construction 17 ()
Let $`S=\{S^1,S^2,\mathrm{},S^L\}`$ be a secret, and let $`\mathrm{\Gamma }^j=\{\stackrel{ˇ}{𝒜}_j,2^𝑽\stackrel{ˇ}{𝒜}_j\}`$, $`j=1,2,\mathrm{},L`$, be the perfect access structures determined from a given access structure $`\mathrm{\Gamma }^R`$. Since each $`\mathrm{\Gamma }^j`$ is a perfect access structure satisfying the monotonicity (3) and (4), we can construct a SSS with $`\mathrm{\Gamma }^j`$ for secret $`S^j`$. Letting $`\{V_1^j,V_2^j,\mathrm{},V_n^j\}`$ be the shares for $`S^j`$ and $`\mathrm{\Gamma }^j`$, the share $`V_i=\{V_i^1,V_i^2,\mathrm{},V_i^L\}`$ realizes the access structure $`\mathrm{\Gamma }^R`$. For $`\mathrm{\Gamma }^R`$, a ramp SSS can also be constructed from $`\{2^𝑽\widehat{𝒜}_j,\widehat{𝒜}_j\}`$ instead of $`\mathrm{\Gamma }^j=\{2^𝑽\stackrel{ˇ}{𝒜}_j,\stackrel{ˇ}{𝒜}_j\}`$.
$`\mathrm{}`$
###### Remark 18
Note that in Construction 17, we have $`\rho _i1`$ for any access structure. For example, in the case that Construction 17 is applied to the $`(k,L,n)`$-threshold access structure, the constructed ramp SSS has $`\rho _i=1`$ although the $`(k,L,n)`$-threshold SSS can be realized with $`\rho _i=1/L`$. Therefore, Construction 17 is not efficient generally. $`\mathrm{}`$
###### Example 19
Consider the following ramp access structure $`\mathrm{\Gamma }_4^R`$ for $`𝑽=\{V_1,V_2,V_3,V_4\}`$:
$`𝒜_3`$ $`=`$ $`\{\{V_1,V_2,V_3,V_4\}\},`$ (107)
$`𝒜_2`$ $`=`$ $`\{\{V_1,V_2,V_3\},\{V_1,V_3,V_4\}\},`$ (108)
$`𝒜_1`$ $`=`$ $`\{\{V_1,V_2,V_4\},\{V_2,V_3,V_4\}\},`$ (109)
$`𝒜_0`$ $`=`$ $`\{𝑨:0|𝑨|2\}.`$ (110)
First, we derive the access structures $`\mathrm{\Gamma }^1`$, $`\mathrm{\Gamma }^2`$, and $`\mathrm{\Gamma }^3`$ based on (107)–(110), and it is easy to see that $`\mathrm{\Gamma }^1`$ and $`\mathrm{\Gamma }^3`$ become $`(3,4)`$\- and $`(4,4)`$-threshold access structures, respectively. Hence, we have $`V_i^1=W_i^{(3)}`$ and $`V_i^3=W_i^{(4)}`$ for $`i=1,2,3,4`$ where $`\{W_i^{(3)}\}_{i=1}^4`$ and $`\{W_i^{(4)}\}_{i=1}^4`$ are the share sets of $`(3,4)`$\- and $`(4,4)`$-threshold access structures for secrets $`S^1`$ and $`S^3`$, respectively. Furthermore, a perfect SSS with the access structure $`\mathrm{\Gamma }^2`$ for a secret $`S^2`$ can be realized by $`\{V_i^2\}_{i=1}^4`$ such that $`V_1^2=W_1^{(3)}`$, $`V_2^2=W_2^{(3)}`$, $`V_3^2=W_3^{(3)}`$, and $`V_4^2=W_2^{(3)}`$ where $`\{W_i^{(3)}\}_{i=1}^3`$ is the share sets of $`(3,3)`$-threshold SSS for $`S^2`$.
According to Construction 17, we can obtain the shares such that $`V_1=\{W_1^{(3)},W_1^{(3)},W_1^{(4)}\}`$, $`V_2=\{W_2^{(3)},W_2^{(3)},`$
$`W_2^{(4)}\}`$, $`V_3=\{W_3^{(3)},W_3^{(3)},W_3^{(4)}\}`$, $`V_4=\{W_4^{(3)},W_4^{(3)},W_4^{(4)}\}`$. Since each share consists of three primitive shares for three secrets $`S^1`$, $`S^2`$, $`S^3`$, the constructed ramp SSS has $`\stackrel{~}{\rho }=\rho ^{}=1`$. $`\mathrm{}`$
The construction of ramp SSSs for general access structures are treated in . But, since the construction in is based on monotone span programming, it is much complicated compared with the multiple assignment map.
### 5.2 Optimal Multiple Assignment Maps for Ramp SSSs
First, let $`𝑾_{(t,L,m)}=\{W_1^{(t,L)},W_2^{(t,L)},\mathrm{},W_m^{(t,L)}\}`$ be the set of primitive shares for the $`(t,L,m)`$-threshold ramp SSS with the coding rate $`\rho _i=1/L`$. Then, defining $`𝒚`$ and $`𝒂(\mathrm{};𝑨)`$ in the same way as the perfect SSSs in Section 3, the optimal ramp SSS by the multiple assignment map for a general access structure $`\mathrm{\Gamma }^R`$ can be obtained by solving the following integer programming problem:
| IP$`{}_{\stackrel{~}{\rho }}{}^{R}\left(\mathrm{\Gamma }^R\right)`$ | | | | | | |
| --- | --- | --- | --- | --- | --- | --- |
| minimize | $`𝒉𝒚^T`$ | | | | | |
| subject to | $`𝒂(1;𝑨)𝒚^T`$ | $``$ | $`0`$ | for | $`𝑨𝒜_L^{}`$ | |
| | $`𝒂(1;𝑨)𝒚^T`$ | $`=`$ | $`j`$ | for | $`𝑨𝒜_j^+𝒜_j^{}`$ | for $`1jL1`$ $`()`$ |
| | $`𝒂(1;𝑨)𝒚^T`$ | $``$ | $`L`$ | for | $`𝑨𝒜_0^+`$ | |
| | $`𝒚`$ | $``$ | 0 | | | |
###### Remark 20
From the monotonicity defined in (103) and (104), it is sufficient to consider only $`𝑨𝒜_j^+𝒜_j^{}`$ instead of all $`𝑨𝒜_j`$ on the marked line $`()`$ in IP$`{}_{\stackrel{~}{\rho }}{}^{R}\left(\mathrm{\Gamma }^R\right)`$. Note that the same primitive shares may be distributed to all shares since there may exist common shares in ramp SSSs. Hence, we may have $`x_N0`$ in the ramp SSSs although we can always assume that $`x_N=0`$ in the perfect SSSs. $`\mathrm{}`$
From Remark 15, significant or common shares $`V_i`$ must satisfy that $`|\phi _\mathrm{\Gamma }(i)|1`$ for any multiple assignment map $`\phi _\mathrm{\Gamma }`$. On the other hand, $`|\stackrel{~}{\phi }_\mathrm{\Gamma }(i^{})|=0`$ must hold for vacuous shares $`V_i^{}`$ for the optimal multiple assignment map $`\stackrel{~}{\phi }_\mathrm{\Gamma }`$ attaining the minimal average coding rate. Hence, it suffices to consider only significant shares and common shares in the ramp SSSs.
###### Example 21
If the access structures $`\mathrm{\Gamma }_4^R`$ in Example 19 is applied to the integer programming problem IP$`{}_{\stackrel{~}{\rho }}{}^{R}\left(\mathrm{\Gamma }_4^R\right)`$, the following multiple assignment map is obtained
$`V_1`$ $`=`$ $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_4^R}(1)=\{W_1^{(7,3)},W_2^{(7,3)}\},`$ (111)
$`V_2`$ $`=`$ $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_3^R}(2)=\{W_3^{(7,3)},W_4^{(7,3)}\},`$ (112)
$`V_3`$ $`=`$ $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_4^R}(3)=\{W_5^{(7,3)},W_6^{(7,3)}\},`$ (113)
$`V_4`$ $`=`$ $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_4^R}(4)=\{W_3^{(7,3)},W_7^{(7,3)}\},`$ (114)
where $`W_i^{(7,3)}𝑾_{(7,3,7)}`$. $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_4^R}`$ attains that $`\stackrel{~}{\rho }=\rho ^{}=2/3`$. $`\mathrm{}`$
Note that the coding rates less than $`1`$ cannot be achieved by Construction 17. Furthermore, our construction is much simpler compared with the method in . But, unfortunately, the integer programming problem may not have any feasible solutions in the case of ramp SSSs.
###### Example 22
The following access structure $`\mathrm{\Gamma }_5^R`$ cannot be constructed by any multiple assignment map since the corresponding integer programming problem has no feasible solution.
$`𝒜_4^{}`$ $`=`$ $`\{\{V_1,V_2,V_3,V_4\},\{V_1,V_2,V_4,V_5\},\{V_2,V_3,V_4,V_5\}\},`$ (115)
$`𝒜_3`$ $`=`$ $`\{\{V_1,V_2,V_3,V_5\},\{V_1,V_3,V_4,V_5\},\{V_1,V_2,V_3\},\{V_1,V_2,V_4\},\{V_1,V_3,V_4\},`$ (116)
$`\{V_1,V_3,V_5\},\{V_2,V_3,V_4\}\},`$
$`𝒜_2`$ $`=`$ $`\{\{V_1,V_2,V_5\},\{V_1,V_4,V_5\},\{V_2,V_3,V_5\},\{V_2,V_4,V_5\},\{V_3,V_4,V_5\},\{V_1,V_3\},\{V_1,V_5\}\},`$ (117)
$`𝒜_1`$ $`=`$ $`\{\{V_1,V_2\},\{V_2,V_3\},\{V_3,V_4\}\},`$ (118)
$`𝒜_0^+`$ $`=`$ $`\{\{V_1,V_4\},\{V_2,V_5\},\{V_3,V_5\}\},`$ (119)
$`\mathrm{}`$
In this case, we can modify the definition of the ramp SSS given by (99) as follows.
$`H(S|𝑨)`$ $`=`$ $`0,\text{for all }𝑨𝒜_L,`$ (120)
$`H(S|𝑨)`$ $``$ $`{\displaystyle \frac{Lj}{L}}H(S),\text{for all }𝑨𝒜_j,1jL1,`$ (121)
$`H(S|𝑨)`$ $`=`$ $`H(S),\text{for all}𝑨𝒜_0.`$ (122)
In order to implement (120)–(122) in the integer programming, it suffices to replace the marked line $`()`$ in IP$`{}_{\stackrel{~}{\rho }}{}^{R}\left(\mathrm{\Gamma }^R\right)`$ by $`𝒂(1;𝑨_j)𝒚^Tj`$. Letting IP$`{}_{\stackrel{~}{\rho }}{}^{R2}\left(\mathrm{\Gamma }^R\right)`$ be the modified integer programming problem, the next theorem holds.
###### Theorem 23
The integer programming problem IP$`{}_{\stackrel{~}{\rho }}{}^{R2}\left(\mathrm{\Gamma }^R\right)`$ always has a feasible solution for any access structure $`\mathrm{\Gamma }^R`$. $`\mathrm{}`$
Proof of Theorem 23: Let $`𝒱`$ be a multiset in $`2^𝑽`$, some elements of which may be the same. Then, for $`𝒱`$ and $`𝑨𝑽`$, we define $`N(𝒱,𝑨)`$ as follows.
$`N(𝒱,𝑨)=\left|\{𝑨^{}𝒱:𝑨𝑨^{}\}\right|,`$ (123)
where all $`𝑨^{}𝒱`$ are treated as different sets even if some of them are the same. Now we construct a multiset $`𝒰`$ for $`\mathrm{\Gamma }^R=\{𝒜_0,𝒜_1,\mathrm{},𝒜_L\}`$ by the next construction.
###### Construction 24
* Let $`𝒰:=\mathrm{}`$ and $`j:=1`$.
* For each $`𝑨𝒜_{Lj}^+`$ satisfying $`N(𝒰,𝑨)<j`$, we add $`𝑨`$ into $`𝒰`$, $`(jN(𝒰,𝑨))`$ times.
* Let $`j:=j+1`$.
* If $`j<L`$, go to (2). In case of $`j=L`$, go to (5).
* Output $`𝒰`$. $`\mathrm{}`$
From the monotonicity of $`\stackrel{ˇ}{𝒜}_j`$ in (103), the family $`𝒰`$ can always be constructed. Then, letting $`𝒰=\{𝑭_1,𝑭_2,\mathrm{},𝑭_m\}`$, we can define a map $`\stackrel{ˇ}{\psi }:\{1,2,\mathrm{},n\}2^{𝑾_{(m,L,m)}}`$ by
$`\stackrel{ˇ}{\psi }(i)={\displaystyle \underset{j:V_i𝑭_j}{}}\left\{W_j^{(m,L)}\right\},`$ (124)
where $`W_j^{(m,L)}𝑾_{(m,L,m)}`$. Note that in the case of $`L=1`$, (124) coincides with the cumulative map in (17). Furthermore, for any set $`𝑭_{\mathrm{}}𝒰`$, we can check from (124) that
$`W_{\mathrm{}^{}}^{(m,L)}{\displaystyle \underset{i:V_i𝑭_{\mathrm{}}}{}}\stackrel{ˇ}{\psi }(i),`$ (125)
holds for all $`\mathrm{}^{}`$ satisfying $`𝑭_{\mathrm{}}𝑭_{\mathrm{}^{}}`$.
Now, assume that $`𝑭_{\mathrm{}}𝒜_j^+`$. Then, from Construction 24, there exist a family of $`j`$ subsets $`\{𝑭_\mathrm{}_1,𝑭_\mathrm{}_2,\mathrm{},𝑭_\mathrm{}_j\}𝒰`$ satisfying $`𝑭_{\mathrm{}}𝑭_{\mathrm{}^{}}`$ for $`\mathrm{}^{}\{\mathrm{}_1,\mathrm{}_2,\mathrm{},\mathrm{}_j\}`$. Hence, it holds from (125) that $`W_{\mathrm{}^{}}^{(m,L)}_{i:V_i𝑭_{\mathrm{}}}\stackrel{ˇ}{\psi }(i)`$ for $`\mathrm{}^{}\{\mathrm{}_1,\mathrm{}_2,\mathrm{},\mathrm{}_j\}`$. This means that we can verify that $`\left|_{i:V_i𝑭_{\mathrm{}}}\stackrel{ˇ}{\psi }(j)\right|mj`$, and $`V_i=\stackrel{ˇ}{\psi }(i)`$ satisfies (120)–(122). Therefore, IP$`{}_{\stackrel{~}{\rho }}{}^{R2}\left(\mathrm{\Gamma }^R\right)`$ always has at least one feasible solution. $`\mathrm{}`$
Note that as shown in the following example, Construction 24 gives inefficient assignments of the primitive shares, generally.
###### Example 25
Assume that the access structure $`\mathrm{\Gamma }_5^R`$ in (115)–(119) satisfies the conditions (120)–(122). First, we apply Construction 24 to the access structure $`\mathrm{\Gamma }_5^R`$. Then, we obtain the following multiset $`𝒰_{\mathrm{\Gamma }_5^R}`$.
$`𝒰_{\mathrm{\Gamma }_5^R}`$ $`=`$ $`\{\{V_1,V_2,V_3,V_5\},\{V_1,V_3,V_4,V_5\},\{V_1,V_2,V_4\},\{V_1,V_2,V_5\},\{V_1,V_4,V_5\},\{V_2,V_3,V_5\},`$ (126)
$`\{V_2,V_3,V_4\},\{V_2,V_4,V_5\},\{V_2,V_4,V_5\},\{V_3,V_4,V_5\},\{V_1,V_4\}\}.`$
Hence, we can obtain $`V_i=\stackrel{ˇ}{\psi }(i)`$, $`i=1,2,\mathrm{},5`$, as follows:
$`V_1`$ $`=`$ $`\stackrel{ˇ}{\psi }(1)=\{W_6^{(11,4)},W_7^{(11,4)},W_8^{(11,4)},W_9^{(11,4)},W_{10}^{(11,4)}\},`$ (127)
$`V_2`$ $`=`$ $`\stackrel{ˇ}{\psi }(2)=\{W_2^{(11,4)},W_5^{(11,4)},W_{10}^{(11,4)},W_{11}^{(11,4)}\},`$ (128)
$`V_3`$ $`=`$ $`\stackrel{ˇ}{\psi }(3)=\{W_3^{(11,4)},W_4^{(11,4)},W_5^{(11,4)},W_8^{(11,4)},W_9^{(11,4)},W_{11}^{(11,4)}\},`$ (129)
$`V_4`$ $`=`$ $`\stackrel{ˇ}{\psi }(4)=\{W_1^{(11,4)},W_4^{(11,4)},W_6^{(11,4)}\},`$ (130)
$`V_5`$ $`=`$ $`\stackrel{ˇ}{\psi }(5)=\{W_3^{(11,4)},W_7^{(11,4)},W_{11}^{(11,4)}\},`$ (131)
where $`W_i𝑾_{(11,4,11)}`$. In this case, we have $`\stackrel{~}{\rho }=21/20`$ and $`\rho ^{}=3/2`$ since it holds that $`H(W_i^{(11,4)})=H(S)/4`$ for each $`i`$.
On the other hand, we can construct the following optimal multiple assignment map $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_5^R}`$ by solving the integer programming problem $`\mathrm{IP}_{\stackrel{~}{\rho }}^{R2}(\mathrm{\Gamma }_5^R)`$.
$`V_1`$ $`=`$ $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_5^R}(1)=\{W_1^{(8,4)},W_2^{(8,4)}\},`$ (132)
$`V_2`$ $`=`$ $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_5^R}(2)=\{W_3^{(8,4)},W_4^{(8,4)},W_5^{(8,4)}\},`$ (133)
$`V_3`$ $`=`$ $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_5^R}(3)=\{W_2^{(8,4)},W_6^{(8,4)}\},`$ (134)
$`V_4`$ $`=`$ $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_5^R}(4)=\{W_7^{(8,4)},W_8^{(8,4)}\},`$ (135)
$`V_5`$ $`=`$ $`\stackrel{~}{\phi }_{\mathrm{\Gamma }_5^R}(5)=\left\{W_9^{(8,4)}\right\},`$ (136)
where $`W_i^{(8,4)}𝑾_{(8,4,9)}`$, and it holds that $`\stackrel{~}{\rho }=1/2`$ and $`\rho ^{}=3/4`$, which are more efficient than the rates of Construction 24. Note that (127)–(131) and (132)–(136) do not satisfy (99) but satisfy (120)–(122). For instance, in (132)–(136), it holds for $`\{V_1,V_5\}𝒜_2`$ that $`H(S|\{V_1,V_5\})=H(S)>H(S)/2`$.
Finally, we compare Construction 17 with Construction 24 for the access structure $`\mathrm{\Gamma }_5^R`$. If we use the cumulative map to realize each perfect SSS with the access structure $`\mathrm{\Gamma }_5^j`$, $`j=1,2,3,4`$, in Construction 17, we obtain $`\stackrel{~}{\rho }=9/5`$ and $`\rho ^{}=2`$. Hence, Construction 17 is more inefficient than Construction 24 in this case. $`\mathrm{}`$
## 6 Conclusion
We proposed a method to construct SSSs for any given general access structures based on $`(t,m)`$-threshold SSSs and integer programming. The proposed method can attain the optimal average and/or worst coding rates in the sense of multiple assignment maps. Hence, the proposed method can attain smaller coding rates compared with the cumulative maps and the modified cumulative maps. Furthermore, the proposed method can be applied to incomplete and/or ramp access structures in addition to complete and perfect access structures.
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# Theoretical Description of the 𝑝𝑑→𝑝𝑑𝜂 Reaction Near Threshold
## 1 Introduction
The great interest in the production of heavy mesons near threshold started with measurements of the $`pd^3`$He$`\eta `$ reaction, which showed a surprisingly strong cross section with an anomalous energy dependence Berger . An impulse approximation description of the process, where the production takes place in nucleon–nucleon scattering, with the other nucleon in the deuteron being essentially a spectator, greatly underestimates the observed rate because of the high momentum components required in the nuclear wave functions Germond . To share the large momentum transfer between the nucleons, a mechanism was proposed whereby a pion was produced in an initial $`ppd\pi ^+`$ reaction to be followed by the production of the observed meson through a secondary $`\pi ^+n\eta p`$ process Kilian . Near threshold, the kinematics are favourable for the final proton and deuteron to *stick* to form the observed <sup>3</sup>He. A quantum–mechanical evaluation of this suggestion FW led to a cross section that was only about a factor of two lower than a precise measurement of the reaction rate Mayer . To understand the situation further, it would be helpful to look at other related final states.
The total cross section for the $`pdpd\eta `$ reaction, where a final <sup>3</sup>He is not formed, was measured at two energies very close to threshold at Saclay Hibou . A preliminary evaluation of the two–step model with an intermediate pion, that was successful in the description of the <sup>3</sup>He$`\eta `$ final state, gave very promising results when compared to these data Ulla . Differential as well as total cross sections for this reaction at higher energies have recently become available from Uppsala Jozef2 , to complement those obtained by the same group for the $`pd^3`$He$`\eta `$ reaction Jozef1 . The time therefore seems opportune to make a further theoretical investigation of the $`pd\eta `$ final state.
The kinematics of the Saclay experiment Hibou are very close to those where a <sup>3</sup>He emerges and so it is not surprising that in this region the two–step model describes the process well, as it did for the <sup>3</sup>He$`\eta `$ final state FW . On the other hand, as one approaches the threshold for $`\eta `$ production in free nucleon–nucleon collisions, one expects a pick–up diagram, corresponding to a quasi–free $`pnd\eta `$ production on the neutron in the target, to dominate. Between these two extremes there is also the possibility of a contribution from the impulse approximation diagram with a quasi–free $`pNpN\eta `$ on a proton or neutron that is bound both initially and finally in the deuteron. Away from threshold, these terms may not be universally suppressed, as they are for <sup>3</sup>He$`\eta `$, because of the greater flexibility in the $`pd\eta `$ kinematics; in certain parts of phase space, the momentum transfers between initial and final deuterons are minimised.
The three types of contributions are described further in sect. 2, where the general kinematics are discussed. The following three sections are devoted to the evaluation of the cross section distributions for the individual terms. Uncertainties in the phases of the amplitudes leads us to neglecting interferences. This may not be too dangerous because the various models tend to populate different parts of what is a three–body phase space. More worrying is that we neglect also all final–state interactions *fsi*. Now, unlike the $`\eta ^3`$He case Mayer , there is no simple prescription for the inclusion of such effects when there is an *fsi* between more than one pair of particles. Experimentally a threshold enhancement is seen only in the $`\eta d`$ invariant mass distribution, whereas that for $`pd`$ looks like phase space, despite the presence of the strong interaction which could lead to the formation of the <sup>3</sup>He Jozef2 . The results for the total cross sections, and both angular and invariant mass distributions, are compared with experiment in sect. 6. The overall production rate is underestimated by about a factor of two, which is rather similar to that found for the $`pd^3`$He$`\eta `$ reaction within a similar theoretical approach. However, the angular distributions of both the proton and deuteron show discrepancies with respect to the data, with the latter showing no signs of the forward deuteron and backward proton peaks corresponding to the spectator proton of the pick–up mechanism. Our conclusions are drawn in sect. 7.
## 2 Defining the problem
The three classes of diagram that we evaluate for the $`pdpd\eta `$ reaction are illustrated in Fig. 1. It is generally assumed that the neutron–exchange diagram (a) dominates the reaction above the free $`NN`$ threshold and, indeed, it is this hypothesis that is the basis of the extraction of the quasi–free $`pnd\eta `$ cross section on a moving neutron from $`pdpd\eta `$ data Stina . Below threshold, the presence of a spectator proton in this case will bias the distributions to low $`d\eta `$ invariant masses.
The impulse approximation diagram of Fig. 1b is often used to model coherent reactions on the deuteron at high energies, where the momentum transfer to the final deuteron can be small. There is, of course, a second term where the $`\eta `$ production takes place on the neutron and some theoretical model is required to deduce the relative phases of the spin–isospin input amplitudes needed to make reliable estimates here FW2 . A final–state interaction between the proton and deuteron to form an <sup>3</sup>He for either the diagrams in Fig. 1a,b would lead to the triangle graph, whose contribution has been shown to be small for $`pd^3`$He$`\eta `$ Germond . That they are not necessarily negligible here is due to the fact that the final $`dp`$ pair does not have to emerge with low excitation energy.
The differential cross section for the $`pdpd\eta `$ reaction is determined by the matrix element $``$ through
$$\text{d}\sigma =\frac{1}{4p\sqrt{s}}\{\frac{1}{6}\underset{\text{spins}}{}\left|\right|^2\}(2\pi )^4\delta ^4(P_iP_f)\underset{j=1}{\overset{3}{}}\frac{\text{d}^3p_j}{(2\pi )^32E_j},$$
(2.1)
where the sum is over the spin projections in the initial and final states. Here $`P_i^2=P_f^2=s`$ is the square of the total cm energy and $`(E_j,𝒑_j)`$ are the energy and three–momentum of the reaction products. The incident flux factor in the centre of mass involves the initial proton cm momentum $`𝒑`$.
Experimental distributions have been presented in terms of the angles of the final particles and the invariant masses of pairs of such particles Jozef2 . For this purpose it is convenient to reduce the phase–space factors to:
$$\text{d}\sigma =\frac{1}{4ps}\{\frac{1}{6}\underset{\text{spins}}{}\left|\right|^2\}\frac{1}{64\pi ^5}p_d\text{d}\mathrm{\Omega }_dp_\eta ^{}\text{d}\mathrm{\Omega }_\eta ^{}\text{d}m_{\eta p},$$
(2.2)
where the final deuteron energy and momentum $`(E_d,𝒑_d)`$, as well as the angles $`\mathrm{\Omega }_d`$, are evaluated in the overall cm system while those of the $`\eta `$ are evaluated in the $`\eta p`$ rest frame, where the invariant mass is given by
$$m_{\eta p}^2=s+m_d^22\sqrt{s}E_d.$$
(2.3)
Analogous formulae follow immediately for the other two combinations of variables.
## 3 The pick–up contribution
Though there is an enhancement at threshold, the energy dependence of the $`pnd\eta `$ total cross section is broadly consistent with $`s`$–wave production up to an excess energy of at least 60 MeV Stina . At threshold there is only one cm amplitude, which can be written as
$$(pnd\eta )=u_{n_c}^{}\left[G\sqrt{s_{pn}}𝒑_n\mathit{ϵ}_d\right]u_p,$$
(3.1)
where $`𝒑_n`$ is the initial momentum and the cm energy $`\sqrt{s_{pn}}`$ arises from the reduction from a relativistic form. Here $`u_p`$ and $`u_{n_c}`$ are proton and charge–conjugate neutron Pauli spinors respectively.
In terms of the amplitude $`G`$, the total $`pnd\eta `$ cross section is
$$\sigma _T(pnd\eta )=\frac{1}{8\pi }p_\eta ^{}p_p^{}|G|^2,$$
(3.2)
where the $`p^{}`$ are evaluated in the $`pn`$ cm frame. The experimental data are then consistent with a value of $`|G|^2=(0.046\pm 0.011)`$ fm<sup>4</sup> Stina .
The spin–averaged amplitude squared for the pick–up term of Fig. 1a reduces to:
$`\frac{1}{6}|(pdpd\eta )|^2`$
$`=\frac{1}{2}[(2\pi )^32m_d](4E_pp_p)^2|G|^2\left\{\stackrel{~}{\phi }_S(q)^2+\stackrel{~}{\phi }_D(q)^2\right\}`$ (3.3)
Here the deuteron wave functions $`\stackrel{~}{\phi }_S(q)`$ and $`\stackrel{~}{\phi }_D(q)`$ are normalised by
$$_0^{\mathrm{}}q^2\left\{\stackrel{~}{\phi }_S(q)^2+\stackrel{~}{\phi }_D(q)^2\right\}=1.$$
(3.4)
They are evaluated at a momentum–squared of
$$𝒒^2=\frac{m_d^2}{E_{di}^2}(\frac{1}{2}pp_p)^2+𝒑_p^2,$$
(3.5)
where the Lorentz boost has been approximated in the low Fermi momentum limit and $`(p_p,𝒑_p)`$ are the components of the final proton momentum $`𝒑_p`$ parallel and perpendicular to the momentum $`𝒑`$ of the incident deuteron, which has energy $`E_{di}`$.
## 4 The triangle diagram
The triangle diagram of Fig. 1b requires as input the amplitudes for the sub–reactions $`NNNN\eta `$. The threshold cm amplitudes for $`I=1`$ and $`I=0`$ are respectively:
$`_1(NNNN\eta )`$ $`=`$ $`\left[𝒲_{1,s}\eta _f^{}\widehat{p}\mathit{ϵ}_i\right]𝝌_f^{}𝝌_i,`$ (4.1)
$`_0(NNNN\eta )`$ $`=`$ $`\left[𝒲_{0,t}\widehat{p}\mathit{ϵ}_f^{}\eta _i\right]\varphi _f^{}\varphi _i,`$ (4.2)
where $`𝝌`$ and $`\varphi `$ represent the isospin–1 and isospin–0 configurations of the $`NN`$ states, with $`\mathit{ϵ}`$ and $`\eta `$ corresponding to the spin–1 and spin–0 combinations.
The spin–averaged total $`\eta `$ production cross sections are
$`\sigma (pppp\eta )`$ $`=`$ $`{\displaystyle \frac{1}{512\pi ^2p_{NN}s_{NN}}}{\displaystyle \frac{m_N}{\left(1+2m_N/m_\eta \right)^{1/2}}}`$
$`\times Q_{NN\eta }^2|𝒲_{1,s}|^2,`$
$`\sigma (pnpn\eta )`$ $`=`$ $`{\displaystyle \frac{1}{1024\pi ^2p_{NN}s_{NN}}}{\displaystyle \frac{m_N}{\left(1+2m_N/m_\eta \right)^{1/2}}}`$ (4.3)
$`Q_{NN\eta }^2\times \left[|𝒲_{1,s}|^2+|𝒲_{0,t}|^2\right],`$
where $`p_{NN}`$ and $`s_{NN}`$ are the momentum and square of the total energy in the cm system, $`m_N`$ the nucleon mass, and $`Q_{NN\eta }`$ the excess energy in the final system FW2 .
In terms of an initial deuteron momentum $`𝒑`$, a final deuteron momentum $`𝒑_d`$, and a momentum transfer $`𝑸=𝒑_d+𝒑`$, the triangle graph is evaluated in the standard way by putting the spectator nucleon on–shell at an average energy $`E_N(\frac{1}{2}𝒑)`$. The spin–averaged matrix element becomes
$`{\displaystyle \frac{1}{6}}{\displaystyle \underset{\text{s}pin}{}}(pdpd\eta )^2={\displaystyle \frac{1}{64}}\left({\displaystyle \frac{m_d}{4E_N(𝒑/2)}}\right)^2\times `$
$`[|3W_{1,s}+W_{0,t}|^2[S_S^2(\frac{1}{2}Q)+S_Q^2(\frac{1}{2}Q)+\frac{2}{3}S_A^2(\frac{1}{2}Q)]`$
$`+|3W_{1,s}W_{0,t}|^2\frac{4}{3}S_M^2(\frac{1}{2}Q)].`$ (4.4)
The spherical, quadrupole, magnetic and axial form factors appearing here are defined by
$`S_S(Q)`$ $`=`$ $`S_a(Q)+S_b(Q),`$
$`S_Q(Q)`$ $`=`$ $`2S_c(Q)S_d(Q)/\sqrt{2}`$
$`S_M(Q)`$ $`=`$ $`S_a(Q)\frac{1}{2}S_b(Q)+S_c(Q)/\sqrt{2}+\frac{1}{2}S_d(Q),`$
$`S_A(Q)`$ $`=`$ $`S_a(Q)\frac{1}{2}S_b(Q)\sqrt{2}S_c(Q)S_d(Q),`$ (4.5)
where the terms on the right hand side represent integrals over the reduced deuteron wave functions $`u(r)`$ and $`w(r)`$:
$`S_a(Q)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}j_0(Qr)\left[u(r)\right]^2\text{d}r,`$
$`S_b(Q)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}j_0(Qr)\left[w(r)\right]^2\text{d}r,`$
$`S_c(Q)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}j_2(Qr)u(r)w(r)\text{d}r,`$
$`S_d(Q)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}j_2(Qr)\left[w(r)\right]^2\text{d}r.`$ (4.6)
Experimental data on $`\eta `$ production in $`pp`$ Calen and $`pn`$ Calen3 collisions above the *fsi* regions show that $`|W_{1,s}|^21.9\times 10^8\mu `$b and $`|W_{0,t}|^223\times 10^8\mu `$b, though the relative phase of these two amplitudes is not determined. In a meson-exchange approach, it is suggested that $`\rho `$–exchange is the largest term FW2 . If further we take the two amplitudes to be in phase then we find that
$`|3𝒲_{1,s}+𝒲_{0,t}|^2`$ $`=`$ $`80\times 10^8\mu \text{b},`$
$`|3𝒲_{1,s}𝒲_{0,t}|^2`$ $`=`$ $`0.4\times 10^8\mu \text{b}.`$ (4.7)
## 5 The two–step model
The contribution of the two–step model with an intermediate pion can be estimated using the same techniques as for the more complicated $`pd^3`$He$`\eta `$ reaction, where there are two integration loops over the unobserved Fermi momenta FW . In terms of the internal momenta marked in Fig. 1c, the $`pdpd\eta `$ matrix element becomes
$`={\displaystyle \frac{\sqrt{2m_d}}{(2\pi )^{3/2}}}{\displaystyle }\text{d}^3q^{}{\displaystyle \underset{\text{int}}{}}(\pi ^+n\eta p)\times `$
$`{\displaystyle \frac{i}{2E_\pi (\sqrt{s}E_{\text{int}}+i\epsilon )}}(pp\pi ^+d)\stackrel{~}{\mathrm{\Psi }}(𝒒^{}),`$ (5.1)
where the sum runs over the internal spin indices.
The pion propagator between the matrix elements for the production and conversion of the pion in eq. (5.1) has been approximated by its positive energy pole. The difference between the external and internal energies, $`\mathrm{\Delta }E=\sqrt{s}E_{\text{int}}`$, depends upon the Fermi momentum $`𝒒`$ and, following ref. FW , we retain only linear terms in this integration variable:
$`\mathrm{\Delta }E`$ $`=`$ $`\sqrt{s}E_d(𝒑_d)E_n(\frac{1}{2}𝒑+𝒒)E_\pi (\frac{1}{2}𝒑𝒒𝒑_d)`$ (5.2)
$``$ $`\mathrm{\Delta }E_0+𝑽𝒒,`$
where the mean energy defect
$$\mathrm{\Delta }E_0=\sqrt{s}E_d(𝒑_d)E_n(\frac{1}{2}𝒑)E_\pi (\frac{1}{2}𝒑𝒑_d)$$
(5.3)
and the relativistic relative velocity between the pion and neutron
$$𝑽=𝒗_\pi 𝒗_n=\left(\frac{E_n+E_\pi }{2E_nE_\pi }\right)𝒑\frac{1}{E_\pi }𝒑_d$$
(5.4)
do not depend on the proton or $`\eta `$ momenta in the final state.
Both the argument of the deuteron momentum space wave function $`\stackrel{~}{\mathrm{\Psi }}(𝒒^{})`$ and the integration variable are the Fermi momentum $`𝒒^{}`$, Lorentz contracted in the beam direction; $`𝒒_{}^{}=𝒒_{}`$, $`𝒒_{}^{}=𝒒_{}/\gamma `$, with $`\gamma =E_d(p)/m_dE_n(\frac{1}{2}p)/m_n`$.
The linearisation is valid only if relatively small Fermi momenta are required to allow the two steps to proceed almost on–shell, which means that the energy defect should be small. At the threshold for $`\eta `$ production $`\mathrm{\Delta }E_015`$ MeV, which illustrates the *kinematic miracle* first noted for the $`pd^3`$He$`\eta `$ reaction Kilian . We then neglect the dependence of the individual $`pp\pi ^+d`$ and $`\pi ^+n\eta p`$ amplitudes upon the Fermi momentum so that the only place where $`𝒒`$ occurs is in the denominator of the propagator. By defining a Lorentz transformed velocity with components
$$𝑽_{}^{}=𝑽_{}\text{and}𝑽_{}^{}=\gamma 𝑽_{}$$
(5.5)
the whole integrand can be written in terms of $`𝒒^{}`$.
After decomposing the deuteron wave function into its $`\mathrm{}=0,\mathrm{\hspace{0.17em}2}`$ components, we are left with integrals of the form
$`i{\displaystyle \text{d}^3q^{}\frac{\stackrel{~}{\phi }_{\mathrm{}}(q^{})T_{\mathrm{}}(\widehat{𝒒^{}})}{\mathrm{\Delta }E_0+𝑽^{}𝒒^{}+i\epsilon }}={\displaystyle _0^{\mathrm{}}}\text{d}te^{i(\mathrm{\Delta }E_0+i\epsilon )t}`$
$`\times {\displaystyle _0^{\mathrm{}}}q^2\text{d}q^{}\phi _{\mathrm{}}(q^{}){\displaystyle }\text{d}\mathrm{\Omega }_q^{}e^{it𝑽^{}𝒒^{}}T_{\mathrm{}}(\widehat{𝒒^{}})`$
$`=S_{\mathrm{}}(\mathrm{\Delta }E_0,|𝑽^{}|)T_{\mathrm{}}(\widehat{𝑽^{}}),`$ (5.6)
where the $`S`$– and $`D`$–state form factors involve integrals over the *configuration–space* deuteron wave functions
$$S_{\mathrm{}}(\mathrm{\Delta }E_0,|𝑽^{}|)=\frac{(2\pi )^{3/2}}{|𝑽^{}|}_0^{\mathrm{}}\text{d}te^{i\omega t}\phi _{\mathrm{}}(t)$$
(5.7)
with $`\omega =\mathrm{\Delta }E_0/|𝑽^{}|`$.
The $`s`$–wave nature of the $`S_{11}(1535)`$ resonance that dominates low energy $`\eta `$ production in $`\pi ^+n\eta p`$, leads to largely isotropic production at low energies. We can then take the matrix element to be proportional to the corresponding spin–non–flip amplitude
$$(\pi ^+n\eta p)=\frac{4\pi \sqrt{s_{\eta p}}}{m}f(\pi ^+n\eta p).$$
(5.8)
For simplicity, of the six invariant $`pp\pi ^+d`$ amplitudes we shall keep only the largest in our energy region, for which
$$(pp\pi ^+d)=\sqrt{2}𝒜(\mathit{ϵ}_d^{}\widehat{𝒑}_\pi )\varphi _{pp},$$
(5.9)
where $`\varphi _{pp}`$ and $`\mathit{ϵ}_d`$ represent the spin–zero and spin–one initial and final $`NN`$ states. This is then related to the cm cross section through
$$\left|𝒜\right|^2=8(2\pi )^2\left\{\frac{s_{pp}}{m^2}\frac{p_p}{p_\pi }\frac{\text{d}\sigma }{\text{d}\mathrm{\Omega }}(pp\pi ^+d)\right\}_{cm}.$$
(5.10)
In addition to the $`\pi ^+`$ diagram of fig. 1c, there is also the possibility of $`\pi ^0`$ propagation between the two steps. This can be taken into account by simply multiplying the matrix element $``$ of eq. (5.1) by an isospin factor of $`3/2`$. Using eq. (2.2), we then arrive at an expression for the distribution of the $`pdpd\eta `$ cross section in terms of the final deuteron angles and $`\eta p`$ invariant mass:
$`{\displaystyle \frac{\text{d}^2\sigma }{\text{d}\mathrm{\Omega }_d\text{d}m_{\eta p}}}={\displaystyle \frac{9}{4(2\pi )^4}}{\displaystyle \frac{s_{\eta p}m_dp_dp_\eta ^{}}{spE_\pi ^2}}\left(\right|S_S|^2+\left|S_D|^2\right)\times `$
$`\left\{{\displaystyle \frac{s_{pp}}{m^2}}{\displaystyle \frac{p_p}{p_\pi }}{\displaystyle \frac{\text{d}\sigma }{\text{d}\mathrm{\Omega }}}(pp\pi ^+d)\right\}_{cm}\left\{{\displaystyle \frac{p_\pi }{p_\eta }}\sigma _{\text{tot}}(\pi ^+n\eta p)\right\}_{cm},`$
where we have taken advantage of the isotropy of the $`\pi ^+n\eta p`$ differential cross section to integrate over the $`\eta `$ angles and to write the result in terms of the total cross section. This option is not open for the other differential distributions, where extra non–trivial integrations have to be performed.
An effective range description of the low–energy $`\pi ^+n\eta p`$ total cross section is
$$\frac{p_\pi ^{}}{p_\eta ^{}}\sigma _{\text{tot}}(\pi ^+n\eta p)\frac{2.76|a|^2}{\left|1iap_\eta ^{}+\frac{1}{2}r_0ap_\eta ^2\right|^2}.$$
(5.12)
The data base has not improved significantly since we took the values $`a=(0.476+0.279i)`$fm and $`r_0=(3.160.13i)`$fm FW .
The $`pp\pi ^+d`$ cross sections are obtainable from the SAID analysis SAID but, because the intermediate pion in Fig. 1a is non–physical, some prescription is needed to extrapolate from the experimental data. We have assumed in the applications that the amplitudes are unchanged when the cm production angle is kept fixed.
## 6 Comparison with experiment
The predictions of the three driving terms of Fig. 1 are compared in Fig. 2 with the existing experimental data Hibou ; Jozef2 . It is immediately apparent that the impulse approximation of the triangle graph in Fig. 1b lies about two orders of magnitude below the experimental data. The ambiguity in the relative phase of the amplitudes $`𝒲_{1,s}`$ and $`𝒲_{0,t}`$ is therefore largely irrelevant and the impulse approximation term will now be dropped from further discussion.
In the near–threshold region of low $`Q`$, the two-step contribution dominates that of the pick–up by an order of magnitude. The kinematics here are somewhat similar to those of low–energy $`pd^3`$He$`\eta `$ reaction where the pick–up term with a proton–deuteron final–state interaction similarly underpredicts the experimental data Laget ; Germond . However, the two–step term levels off at $`Q80`$ MeV and eventually decreases at higher energies due to a combination of factors. The region where the intermediate pion is close to being physical becomes a smaller fraction of the allowed phase space but also the $`ppd\pi ^+`$ cross section falls steeply above the $`\mathrm{\Delta }`$ region. Since the pick–up term must approach that for quasi–free $`pnd\eta `$ above the $`NN`$ threshold, the two contributions must cross before then. In our estimates, this occurs at $`Q95`$ MeV.
The incoherent sum of the three contributions is also shown in Fig. 2. Though the shape is very similar to that of the experimental data, these lie about a factor of two above the curve at high energies and a little bit more near threshold. This parallels a similar underprediction of the $`pd^3`$He$`\eta `$ reaction in a two–step model that includes the effects of the $`\eta ^3`$He *fsi* FW . The fractions of phase space where the $`pd`$ or $`\eta d`$ final–state interactions might enhance the $`pd\eta `$ channel are relatively small at the higher energies.
More information can be derived from the differential distributions that were measured at $`T_p=1032`$ MeV ($`Q=72.3`$ MeV) Jozef2 . The only significant enhancement in the three invariant mass distributions of Fig. 3 is that near the $`\eta d`$ threshold, which is consistent with the large scattering length suggested by the low energy $`pnd\eta `$ total cross section data Stina . In contrast, despite the existence of the nearby <sup>3</sup>He bound–state pole, there is no sign of any $`pd`$ *fsi*. The two–step model leads to only minor distortions of the invariant mass phase spaces, pushing $`m_{p\eta }`$ to slightly larger values, whereas the pick–up contribution preferentially populates low $`m_{d\eta }`$ masses which is reflected kinematically as a slight increase at higher values of $`m_{pd}`$.
Apart from the overall strength being too low by a factor of two, the sums of the two contributions give plausible descriptions of the invariant mass distributions but the same cannot be said for the angular distributions shown in Fig. 4. In the pick–up contribution of Fig. 1a, the final proton is a spectator and the transverse momentum is governed by the Fermi momentum components in the deuteron. This automatically gives a sharp peak for a cm proton angle close to the backward direction and this has a slightly weaker kinematic reflection around the forward deuteron direction. Given that the estimate within the pick–up model has relatively few uncertainties, the discrepancy with the experimental data of Ref. Jozef2 is particularly significant.
There are far more uncertainties in the evaluation of the two–step model since the intermediate pion is generally off its mass shell and this effects the kinematics of the $`ppd\pi ^+`$ reaction. However, the deuteron angular distribution of Fig. 4a does suggest that the models need to be supplemented at large angles and this might be the reason for the underestimation of the total cross section.
## 7 Conclusions
We have estimated the contributions of three different models to the total and partial cross sections of the $`pdpd\eta `$ reaction below the $`\eta `$–production threshold in nucleon–nucleon collisions. The impulse approximation turns out to be largely negligible compared to the other two terms of which the two–step model is shown to be dominant in the near–threshold region. However, although some of our kinematic approximations may start to break down before one reaches the free $`pn`$ threshold ($`Q192`$ MeV), the estimates suggest that the two–step mechanism would provide only a very small correction to the pick–up interpretation of the CELSIUS quasi–free $`pnd\eta `$ data Stina .
We have neglected the final–state interactions which should distort the lower edges of the invariant mass distributions shown in Fig. 3 while having a smaller effect on the total cross section away from the threshold region. In fact the only *fsi* clearly seen in the experimental data is that associated with the $`\eta d`$ channel Jozef2 . Now the total cross section at the higher energies is underestimated by a factor of two, which is a very similar factor to that found for the near–threshold two–body $`pd^3`$He$`\eta `$ reaction interpreted in the same two–step approach, after the inclusion of the necessary *fsi* effects FW . It is, however, intriguing to note that, although there is no sign of the strong $`pd`$ *fsi* in the experimental data of Fig. 3b, the normalisation of the total $`pdpd\eta `$ total cross section is predicted successfully from the $`pd^3`$He$`\eta `$ data using the final–state extrapolation theorem extrapolation without implementing a full dynamical model.
Apart from the *fsi* regions, it is hard to draw firm conclusions on the models from the invariant mass distributions of Fig. 3. There is far more information to be gleaned from the angular distributions of Fig. 4. The models underpredict the data in the backward deuteron hemisphere and, more critically, there is no sign in the data for the sharp peaks for forward–going deuterons and backward–going protons being produced in the pick–up process. Though these might be softened somewhat by multiple scatterings or *fsi*, it should also be noted that if either the deuteron or the proton is lost down the beam pipe then the event is not registered Jozef2 .
If indeed the two–step model is the dominant mechanism well below the nucleon–nucleon threshold, it would have consequences for the $`pdppn\eta `$ four–body final state, where one could expect to see a strong $`pn`$ *fsi*. Data on this reaction were taken at CELSIUS simultaneously with those on the other channels, but their detailed analysis is not yet complete Jozef3 . Data also exist from COSY on the analogous sub–threshold $`pdK^+\mathrm{\Lambda }d`$ reaction Valdau and it would seem likely that the corresponding two–step model should play an important role there as well.
###### Acknowledgements.
We are much indebted to our colleagues in Uppsala, especially J. Złomańczuk, for discussions regarding the $`\eta `$–production programme at CELSIUS. One of the authors (CW) is grateful for financial support and hospitality from both the Department of Radiation Sciences and the The Svedberg Laboratory of the University of Uppsala. Comments by Y. Uzikov regarding the impulse approximation term proved helpful.
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# Size distribution of circumstellar disks in the Trapezium cluster
## 1 Introduction
Understanding the formation of planets from protoplanetary disks surrounding young stars is an important goal of modern day astrophysics. Although there is evidence for most stars in our Galaxy to have formed in dense photo-evaporating clusters much like the Trapezium cluster (e.g., Lada & Lada 2003), it is not yet clear that the circumstellar disks once associated with these stars survived the harsh conditions imposed by massive star formation, and matured into planetary systems. The Trapezium cluster in the Orion star forming region ( M42, NGC 1976) is a unique laboratory to tackle this problem, mainly because it is the closest (d$`=`$450 pc, e.g., Muench et al. 2002) massive star forming region. It is also very young ($``$ 1 Myr) (Hillenbrand 1997; Hillenbrand & Carpenter 2000), and harbors nearly 2 000 YSOs within a few parsecs (Hillenbrand & Hartmann 1998), the cluster core alone containing over 100 YSOs stars within 0.1 pc of the star $`\theta ^1`$ Ori C.
Solar system-sized circumstellar disks in the Orion Nebula were first inferred from radio observations of compact ionized regions surrounding young low-mass stars (Churchwell et al. 1987). The Hubble Space Telescope (HST) subsequently provided the most compelling evidence for disks in the spectacular images of extended circumstellar material surrounding young low-mass stars in the core of the Orion Nebula. Over half of the 300 YSO observed in the HST images were classified as proplyds (PROtoPLanetarY DiskS) (O’Dell et al. 1993, O’Dell & Wen 1994; O’Dell & Wong 1996; McCaughrean & O’Dell 1996; Bally et al. 1998a; Bally et al. 2000), flattened circumstellar clouds of dust and gas surrounding young stars rendered visible through inclusion in or proximity to an H II region. The intense UV radiation fields of the massive OB stars heat the disk surface, drive mass-loss and produce bright ionization fronts due to the emission-line radiation arising from the outer parts of the proplyds. Most of them are comet-shaped ionized envelopes pointing directly away from the brightest OB stars and are believed to contain evaporating circumstellar disks (McCaughrean et al. 1998; Johnstone et al. 1998). The principal sources of ionizing radiation in the region are the O6, 45 M star $`\theta ^1`$ Orionis C, the brightest of the Trapezium compact group of 4 high-mass OB stars, and the O9.5, 25M star $`\theta ^2`$ Orionis A located several arcminutes to the south at a distance of 0.3 pc from $`\theta ^1`$ Ori C. Dusty disks are seen either as dark silhouettes against the bright background nebular light – the pure silhouettes – (McCaughrean & O’Dell 1996; Bally et al. 2000) or embedded in the light from their own ionization fronts – the embedded silhouettes (Bally et al. 2000). Proplyds display a variety of forms. The ones closer to $`\theta ^1`$ Ori C have bright cusps in optical emission lines (H<sub>α</sub>, \[O III\], \[N II\]) facing $`\theta ^1`$ Ori C and “tails” extending from the ends of the cusps. The farthest have curved and close boundaries that include the entire circumstellar cloud. In the fields of the inner portion of the Orion Nebula – a region that contains more than 300 YSO of the nearly 2000 members of the extended Trapezium cluster – imaged with the HST, 161 proplyds were identified, of which 15 are pure silhouettes and 146 are bright droplets surrounded by ionization fronts (IF).
O’Dell 2001a examined WFPC2/HST parallel planned observations in the outer portions of the Orion Nebula and near the center of its companion H II region M43 located several arcminutes north of the Trapezium and powered by a single B0 spectral type star, NU Ori. His discovery of three new bright proplyds (093-822; 307-1807; 332-1605) and one pure sillhouette proplyd (321-602) is an important proof that additional young stars, disks and jets remain to be revealed in the outer parts of the Orion Nebula. Subsequently, Smith et al. smith05 (2005), as a result of an H<sub>α</sub> survey with the HST Wide Field Camera of the Advanced Camera for Surveys (ACS/WFC), discovered 10 new silhouette disks in the outskirts of the Orion Nebula and in its neighboring region M43: 6 pure silhouettes – 053-717, 110-3035, 141-1952, 280-1720, 347-1535, 216-0939 – and 4 embedded silhouettes – 132-042, 124-132, 253-1536, 181-826 (this one discussed in Bally et al. 2005). These new disks exhibit extended emission from bipolar reflection nebulae and microjets due to a fainter background in the less intense radiation fields regions far from the Trapezium core. Up to date publications, 171 proplyds were imaged with the Hubble Space Telescope (HST) in the Orion Nebula and M43: 150 bright cusps and 21 pure silhouettes.
In this paper we analyze available HST-WFPC2 images to measure the diameters of 149 proplyds (14 are pure silhouettes and 135 are bright proplyds surrounded by ionization fronts) and derive the basic statistics of this population. These statistics can provide important constraints to models of protoplanetary disk evolution and planet formation in young stellar clusters. The observations are described in §2 and the results presented in §3. The results are discussed in §4 and the main conclusions are summarized in §5.
## 2 Observations
To perform the diameter measurements we selected from the available Hubble Space Telescope (HST) images of the Orion Nebula (Bally et al. 2000)<sup>1</sup><sup>1</sup>1The HST mosaicked “master” images of the Orion Nebula are accessible at http://casa.colorado.edu/$``$bally/HST/HST/master/. the one that provided the best contrast between the disk and the background nebula light. The available images are mosaic images from data acquired through filters F673N, F658N, F656N, F631N, F547N and F502N, under the general observer (GO) programs, GO 5085 (O’Dell & Wong 1996) and GO 5469 (Bally et al. 1998a; O’Dell 1998). After quantitative comparisons we selected the H<sub>α</sub> image since it provided the best contrast and had the best signal-to-noise. The selected H<sub>α</sub> image covers an area in the sky of 7.8 $`\times `$ 8.3 arcmin<sup>2</sup> and has a resolution of about 0.<sup>′′</sup>1. World Coordinate System (WCS) coordinates were inserted using IRAF<sup>2</sup><sup>2</sup>2IRAF is distributed by the National Optical Astronomy Observatory, which is operated by the Association of Universities for Research in Astronomy, Inc., under cooperative agreement with the National Science Foundation. scripts and the images were searched for circumstellar disks seen only in silhouette, bright proplyds with dark disks, and bright proplyds without visible disks. The names of the objects discussed in this paper use the coordinate-based system introduced by O’Dell & Wen odell94 (1994), where the numbers are truncated positions (in the equinox J2000.0 FK5 reference system) rounded off to 0<sup>s</sup>.1 of right ascension and 1<sup>′′</sup> in declination. The first three coordinates correspond to the right ascension (J2000.0) given to 0<sup>s</sup>.1, with 5<sup>h</sup>35<sup>m</sup> subtracted and the decimal indicator dropped. The last three digits correspond to the declination (J2000.0) with 520 added. For example, 131-046 has a right ascension of 5<sup>h</sup>35<sup>m</sup>13<sup>s</sup>.1 and a declination of –52046<sup>′′</sup>. When the object lies north of –520, a four digit declination indicator is used, like in 121-1925.
For the image analyzed in this paper the typical errors in the positions are about 0.5 pixels or 0.<sup>′′</sup>05. For a distance of 450 pc to the Orion Nebula, the physical resolution of the HST image is 0.<sup>′′</sup>15 or $``$ 67.5 AU. ($``$ 1.5 pixel/FWHM; 1 pixel corresponds to $``$ 0<sup>′′</sup>.1 or 45 AU). Bally et al. bally00 (2000) analyzed dithered images obtained with the WFPC2 in the HST Cycle 6 under program GO 6603. The dithered images were combined into final images using the drizzling technique resulting in individual wide-field frames with 20% better resolution than the non-drizzled image analyzed in this paper. This has an impact on the measurements of the proplyds’ sizes, i.e., the diameters measured in this paper are, on average, 20% greater than the diameters listed in Bally et al. bally00 (2000).
## 3 Results
### 3.1 The sample
We assume circumstellar disks to be essentially circular structures (typical disk thickness, $`10`$% of diameter, is negligible when compared to the precision of our measurements, see section 3.2). Projections of a circular disk at any random orientation always corresponds to an ellipse, exception given to the perfect edge-on case. The diameter of the major axis of such an ellipse is always equal to the diameter of the circular disk seen in projection. Hence, a direct measurement of the major axis of an observed ellipse is a direct measurement of the disk diameter.
The initial sample consisted of 161 sources (see Table 3, available at the CDS), all listed in previous papers (O’Dell et al. 1993, O’Dell & Wen 1994; O’Dell & Wong 1996, McCaughrean & O’Dell 1996; Bally et al. 2000) but only 144 of them were identified in our analysis (14 were not found, 2 were not included and 1 was out of the image). We found an object which was clearly proplyd-like but its coordinates did not correspond to any object in previous catalogs. We named it 266-557 using O’Dell’s coordinate-based system and added it to the sample. For the 14 proplyds not found in our analysis, 2 are silhouettes (132-1832; 172-028) and 12 are bright proplyds and they are, either indistinguishable from the background (064-705; 153-1902; 158-425; 166-519; 171-315; 174-400; 179-536; 205-421; 208-122; 215-317; 221-433), or very close to a bright star (163-322). We suggest that some of their coordinates can be possibly wrong. The proplyds 168-326NW and 168-326SE, listed as two different objects in previous papers, were not resolved in the image as two objects, so they were not included in the sample. From the 15 silhouettes discovered until and listed in the paper Bally et al. bally00 (2000), the silhouette 294-606 appears out of the field of view in the H<sub>α</sub> mosaic image we used and, therefore, was not considered in the sample. The proplyds 132-1832, 172-028, 171-315 and 205-421 were measured directly from their published images in the papers O’Dell & Wen odell94 (1994) and Bally et al. bally00 (2000), resulting in a final sample of 149 proplyds where 14 are pure silhouettes and 135 are bright proplyds. Table 3, available at the CDS, lists the observed properties of all the 162 externally illuminated YSOs considered in this paper, measured coordinates (RA and DEC), disk and ionization front diameters in pixel and AU (measured in H<sub>α</sub> emission), projected distances to $`\theta ^1`$ Ori C and comments. See discussion on the procedure followed in the next section.
O’Dell 2001a and Smith at al. smith05 (2005) proplyds were not included in our sample because they lie farther north from the region covered by our H<sub>α</sub> image.
### 3.2 Procedure
The measuring procedure can be summarized as the following: for the dark silhouettes the diameter in pixel was considered to be the diameter of an ellipse fitted to the $``$ 10% contour below the background intensity. Because of the highly variable background in the Orion Nebula, the considered value for the intensity, is an average (Imax + Imin)/2, of the immediate background surrounding the disk. For the bright proplyds, within the variety of dimensions that can be measured, the most useful and practical is the “chord diameter”, meaning the distance between the tips of the cusps of the proplyds as described by O’Dell odell98 (1998). This can be determined more accurately than any other dimension, such as the size along the direction from the ionizing star to the cusp, and should not vary with different spatial orientations. But, the uncertainty of the chord diameter is primarily determined by the judgment involved in where to define the cusp boundary and amounts 1 to 2 pixels. To overcome this problem, we defined the cusp boundary as a circle fitted to the contour of plus $``$ 30% the average background intensity. The “chord diameter” is then the diameter of this circle with an uncertainty of 0.5 pixel or 22.5 AU, the error associated with the measurement only. Unfortunately, because of the highly variable background intensity, the contrast between proplyds/silhouettes and the background is sometimes low. In our sample of 135 bright proplyds, 37 objects, or 27% of the total sample, were not bright enough so they were fitted to a contour of +10% the average background intensity. At the same time, although clearly detected, 4 silhouettes were not dark enough and were measured to plus a few percent.
### 3.3 Do the proplyd diameters vary with distance from $`\theta ^1`$ Orionis C ?
The observational establishment of the proplyds as a well defined class of objects has led to several attempts to theoretically model them (Henney et al. 1996; Johnstone et al. 1998; Henney & Arthur 1998; Störzer & Hollenbach 1999; Richling & Yorke 2000). In this section we will introduce the most important aspects of proplyd photoevaporation models to stress that a correlation between proplyd diameter and distance from the ionizing source is expected in theory.
The diameter of the hydrogen ionization fronts (bright proplyds) is a complex function dominated by the incident UV photon flux or true distance to the illuminating source and the disk radius. The appearance of the proplyds can be explained by the interaction of EUV – extreme-ultraviolet (Lyman continuum; h$`\nu `$ $`>`$ 13.6 eV) and FUV – far-ultraviolet (6 eV $`<`$ h$`\nu `$ $`<`$ 13.6 eV) radiation from an external source with a circumstellar disk. The FUV photons are absorbed mainly by dust and penetrate much deeper into the disk than the EUV photons. At relative high column densities they dissociate molecules, heat the material impenetrable to the EUV photons and initiate a neutral flow away from the disk. This “photon-dominated region” or PDR is encased within an hydrogen IF where the outflowing PDR material can be ionized by the EUV photons. Depending on the ratio I(UV)= FUV/EUV, the IF can stand off at a considerable distance from the disk surface, appearing as a bright round head in the direction of the illuminating source with an extended “cometary” tail of lower emissivity in the far side of the disk. Johnstone et al. johnstone98 (1998), presented analytical and numerical models of the structure of the neutral flows. Depending on the disk radius, the I(UV)= FUV/EUV ratio and the column of neutral gas in the PDR, the neutral flow can be either FUV-dominated or EUV-dominated. For a given disk outer radius, as the distance to the ionization source increases, the EUV flux declines, I(UV) increases and the radius of the IF is expected to increase. The most rapid disk erosion occurs for large systems close to the Trapezium. Störzer & Hollenbach storzer99 (1999) improved the previous Johnstone et al. johnstone98 (1998) model including the results of both equilibrium and non-equilibrium photo-dissociation region (PDR) codes to calculate the column density and temperature inside the PDR. They determined under which circumstances FUV-dominated flows are possible and explained the observed IF diameters of the Orion proplyds by the FUV-dominated flows. These FUV-dominated flows are extended with IF radius r<sub>IF</sub> $``$ d$`{}_{}{}^{2/3}{}_{\theta ^1OriC}{}^{}`$ and typically, r<sub>IF</sub> $``$ 2r<sub>d</sub>, for a distance of 0.01–0.3 pc from $`\theta ^1`$ Ori C. The mass-loss rate is proportional to the disk radius and not so dependent on the distance to $`\theta ^1`$ Ori C. Outside this region, EUV photons dominate the photoevaporation, the IF is closer to the disk surface, r<sub>IF</sub> $``$ 2r<sub>d</sub>, and the mass-loss rate is proportional to d$`{}_{}{}^{1}{}_{\theta ^1OriC}{}^{}`$ and r$`{}_{}{}^{3/2}{}_{d}{}^{}`$.
Is there any observed systematic variation of the proplyds diameter with the distance to the main ionizing stars?
Figure 1 presents the IF chord diameter as a function of the projected distance to $`\theta ^1`$ Ori C for the 135 bright proplyds sample. The different symbols indicate the main ionizing stars and, for each proplyd, they were determined by the relative orientation of the bright cusps and tails. For the not so clear cases, we drew a line in our image connecting the ionizing stars and the center of the cusp and compared them with the proplyd’s tail-cusp direction. The result of this procedure was that there are 105 proplyds primarily ionized by $`\theta ^1`$ Ori C, 19 by $`\theta ^2`$ Ori A, 9 by both $`\theta ^1`$ Ori C and $`\theta ^2`$ Ori A (like 244-440) and 2 ionized by $`\theta ^1`$ Ori C and $`\theta ^2`$ Ori B (252-457 and 282-458). Fig.1 clearly shows no obvious correlation between the “chord diameter” and the projected distance to the ionizing Trapezium OB stars. Assuming a random distribution of disk sizes across the cluster, this result hints at poor correlation between IF size and distance, unlike theoretical predictions.
### 3.4 Deriving disk diameters from the IF chord diameters
It is not trivial to infer disk diameters from the bright ionization front diameters. Nevertheless, empirically, one can investigate if there is a typical diameter ratio R (R = embedded disk diameter/ionization front diameter) since we have in the sample bright proplyds with clearly resolved embedded disks. We identified 10 such cases (072-135; 141-520; 143-552; 163-222; 176-543; 181-247; 182-413; 197-427; 206-446; 244-440), at different projected distances from the ionizing sources, and measured both the diameter of the ionization front cusp and the diameter of the disk, as described in section 3.2. Figure 2 represents the R values for the 10 silhouettes with embedded disks as a function of the disk diameters. The different symbols are related to the different main ionizing stars. We find the average diameter ratio for these 10 sources to be $`<`$R$`>`$ = 49 $`\pm `$ 16%, a value essentially identical to theoretical expectation for a distance of 0.01–0.3 pc from $`\theta ^1`$ Ori C (Störzer & Hollenbach 1999). The observational parameters for these 10 proplyds are listed in Table 1.
Figure 3 is a spatial distribution diagram of the 10 proplyds with embedded silhouette disks. The black filled circles represent the silhouette disks and the colored, unfilled disks represent the ionization fronts. Green is for the proplyds ionized primarily by $`\theta ^1`$ Ori C (072-135; 163-222; 181-247; 182-413), blue for the proplyds ionized by $`\theta ^2`$ Ori A (176-543; 206-446) and red for proplyds ionized by both stars (141-520; 143-522; 197-427; 244-440). The circle diameters are proportional to the disk and IF diameters.
There is a fairly large dispersion in the ratio R = disk diameter/IF diameter, as showed by Fig. 2, but it is not clear that R is correlated with the projected distance, or true distance, to $`\theta ^1`$ Ori C (see Fig.3). For approximately the same proplyd disk diameter there are examples of IFs very close to the disk and IFs that stand at a considerable distances from the disk. There is also no apparent increase of the IF’s diameters with the projected distance to $`\theta ^1`$ Ori C (Fig.1). The proplyds 244-440 (R = 15%) and 072-135 (R = 71%) illustrate the two opposite “extreme” cases. 244-440 is localized at a distance of 142.<sup>′′</sup>3 from $`\theta ^1`$ Ori C and 29.<sup>′′</sup>2 from $`\theta ^2`$ Ori A and it has the larger IF observed in the Trapezium with a diameter of 2520 AU or 5.<sup>′′</sup>6. One of the reasons could be because it is being ionized by both stars as seen from the shape and orientation of its cusp. Nevertheless, its disk is only 387 AU or 0.<sup>′′</sup>86. We can find enough cases of proplyds that lie in projection far from $`\theta ^1`$ Ori C and still have small IF, despite the large diameter of their disks. For example, 072-135 is positioned at a distance of 176.<sup>′′</sup>3 or 0.38 pc in the north-west direction from $`\theta ^1`$ Ori C and has a small IF of 1.<sup>′′</sup>05 with a large disk of 0.<sup>′′</sup>75. At the same time, 176-543, approximately at the same distance from $`\theta ^1`$ Ori C as 244-440 or 141.<sup>′′</sup>3, has an IF of only 1.<sup>′′</sup>17 or $``$ 1/5 of 244-440’s chord diameter. Its disk is 0.<sup>′′</sup>67 and R = 57%. In summary, there is no obvious indication for larger proplyds to be located farther from the UV source, even taking into consideration the unknown distant projection correction. In conclusion, it seems reasonable to adopt the average value of R, $`<`$R$`>`$ = 49 $`\pm `$ 16%, from our sample of 10 proplyds distributed at different distances, and use it as a calibrator to compute the disks diameters for the 125 bright proplyds without embedded silhouettes. We find the same general result in the data from previous measurements by O’Dell odell98 (1998) and Bally et al. 1998a , discussed in the next section.
### 3.5 Comparison with previous work
In this section we make a detailed comparison of the data presented in this paper and the observational data published in previous papers – Bally et al. 1998a , O’Dell odell98 (1998), Johnstone et al. johnstone98 (1998), Störzer & Hollenbach storzer99 (1999) and Bally et al. bally00 (2000) (see Table 2).
For the subsample of proplyds listed in these papers, and common to our sample, an identical analysis was performed. Because the assumed distances to the clusters used in the different papers are not always the same (we use d<sub>Orion</sub> = 450 pc through out this paper) we scaled all results to the same distance and only then compared them. To compare the different subsamples of proplyds with embedded disks we proceeded as follows: we computed the diameter ratio R for each proplyd, the average diameter ratio $`<`$R$`>`$ for all the subsample of proplyds, and the linear correlation coefficient r between the disk and IF diameters and the projected distances to the main ionizing stars<sup>3</sup><sup>3</sup>3The correlation coefficient r is a measure of the linear association between variables. The strength of this association is sometimes expressed as the square of the correlation coefficient. The resulting statistic is known as variance explained. For example, a correlation of 0.5 means that (0.5)<sup>2</sup> or 25% of the variance in Y is “explained” or predicted by the X variable. This parameter is used throughout this paper..
Bally et al. 1998a observed four fields near the Trapezium with the Planetary Camera (PC), the 0.<sup>′′</sup>05 angular resolution portion of the WFPC2/HST, on March 1995. Observations were obtained through narrow-band filters centered on the \[SII\], \[N II\], H<sub>α</sub>, \[O I\] and \[O III\] lines and in an emission-line-free continuum channel. They present the dimensions of the various emission and opaque components, and the peak surface brightness of the objects in each emission line, for 43 Trapezium proplyds, of which 40 are bright (21 with dark disks embedded) and 3 are pure silhouettes. From these, 38 are common to our sample (35 bright proplyds + 3 pure silhouettes) and 19 are bright proplyds with dark embedded disks. For the maximum projected distance to $`\theta ^1`$ Ori C in the sample (184-427 at 70<sup>′′</sup>), the average diameter ratio $`<`$R$`>`$ computed for the 19 proplyds is 43 $`\pm `$ 18%. Bally et al. 1998a claims that, when both silhouette and IF are seen, the ratio of the semi-major axis of the silhouettes divided by the IF radius ranges from 0.25 to 0.67. This is also confirmed in our analysis. The values of R show a large dispersion but are in that range, with exceptions given to 158-327, 160-353 and 166-316. The authors argue that there is a loose correlation between the cusp radius and the projected distance to $`\theta ^1`$ Ori C, with the cusp radius increasing as r(d)$``$ d<sup>α</sup> with $`\alpha `$ = 0.5–0.8, where d is the projected distance from $`\theta ^1`$ Ori C. For the 34 bright proplyds common to our sample, we find $`\alpha `$ = 0.46. The linear correlation coefficient between the two variables is 0.54 or $``$ 30%, meaning a loose correlation for this specific sample of proplyds.
O’Dell odell98 (1998) presents a set of 22 bright proplyds observed in the Bally’s GO program 5469 with the HST/PC in 1995, of which 20 are common to our sample and 19 to Bally et al. 1998a . The author found a very loose correlation of the proplyd cusp diameter with the distance from $`\theta ^1`$ Ori C, for a maximum projected distance in his sample of 70<sup>′′</sup> (184-427) or $``$ 0.15 pc. We computed the linear correlation coefficient for the 20 common sources and found a value of 0.47 or 22%, indicator of a minor correlation between the variables.
Johnstone et al. johnstone98 (1998) compare their model results with HST observations from Bally et al. 1998a . For 15 proplyds disks, for a maximum projected distance to $`\theta ^1`$ Ori C of 70<sup>′′</sup> (184-427), the average diameter ratio $`<`$R$`>`$ equals 48 $`\pm `$ 22%. The correlation coefficient between disk diameters and projected distance to $`\theta ^1`$ Ori C is 0.26 or 7%. This paper presents a sample of 41 proplyds of which 30 (27 common to our sample) were measured with the HST Planetary Camera and 11 (3 common to our sample) with the Wide Field Camera. The total 30 bright proplyds IF diameters are correlated with the distance to $`\theta ^1`$ Ori C by 42% or r = 0.64, indicating a large correlation between them.
Störzer & Hollenbach storzer99 (1999) used the best 10 measured disks reported by Johnstone et al. johnstone98 (1998) , in which both disk and IF diameters could be observed (Bally et al. 1998a). For a maximum projected distance of 56.<sup>′′</sup>7 (182-413 or HST10), $`<`$R$`>`$ = 48 $`\pm `$ 20%. There is a very high correlation between both disk and IF diameters and the projected distance to $`\theta ^1`$ Ori C, r<sub>IF</sub> = 55% and r<sub>disk</sub> = 70%. But, we have to keep in mind that this is a very particular sample.
The last test was performed in the data from Bally et al. bally00 (2000). Although, the only diameters listed are for the 15 pure silhouettes, they published images of all the embedded disks seen against the bright ionization fronts: 24 disks in H$`\alpha `$, 16 disks in \[O III\] and 2 disks in \[N II\]. The IF, and disk diameters of 17 proplyds (disks boundaries are not clear in all pictures) were measured directly from the published H$`\alpha `$ images. The proplyd 141-520 measures were taken from the \[N II\] image, increasing the sample, from Bally et al. bally00 (2000) paper, to 18 bright proplyds. There is no significant correlation between the disk and IF diameters and the projected distance to $`\theta ^1`$ Ori C. For this sample: r<sub>IF</sub> = 2% and r<sub>disk</sub> = 8%, for a maximum projected distance of 176.<sup>′′</sup>3 (072-135) or 0.38 pc. The average diameter ratio is $`<`$R$`>`$ = 44 $`\pm `$ 17% and the IF diameters are, on average, 1.2 $`\pm `$ 0.2 times smaller than the ones in this paper, in agreement with the 20% increased resolution of Bally et al. (2000) drizzled images. For our sample of 135 bright proplyds, and separating them in subsamples accordingly with the main ionizing O star, we computed the linear correlation coefficient between the disk and IF diameters and the projected distance to the main ionizing source. The first subsample is composed of 105 bright proplyds which are being ionized primarily by $`\theta ^1`$ Ori C, and we found that there is no correlation between the two variables for a maximum projected distance of 264<sup>′′</sup> (005-514) (Figure 1). When we consider just the proplyds which are at a distance less than 0.3 pc (Störzer & Hollenbach 1999), the correlation is low (-0.26 or 7%). For the 19 proplyds ionized by $`\theta ^2`$ Ori A, for a maximum projected distance of 151<sup>′′</sup> (224-728), we get a week correlation of 13% for IF diameter with the projected distance to $`\theta ^2`$ Ori A, even for distances less than 0.3 pc. The disks show a higher correlation of 18%.
In summary, a dispersion of R for the Trapezium proplyds is found by all the authors for the previous papers and, in all of them as well, there is no obvious correlation between the IF “chord diameter” and the projected distance to the Trapezium O stars. Moreover, and within the errors, all previous papers agree on the value of $`<`$R$`>`$. It seems reasonable to assume an average value for the diameter ratio disk diameter/ IF diameter and use it as a “calibrator” to infer the disks diameters for the 125 bright proplyds without visible disks. The rms in $`<`$R$`>`$ will be used to estimate the error in this approximation.
### 3.6 Caveats
Selection effects make the detection of the disks difficult. Dark silhouettes can be easily overwhelmed by the unknown line-of-sight contribution from the nebula or from young stars at their centers. This is especially true for small ($`<`$ 2 WFPC2 pixels, or 90 AU) disks. On the other hand, small high surface brightness proplyds near $`\theta ^1`$ Ori C are hard to distinguish from point sources. In general, the location of a proplyd with respect to the illuminating stars strongly influences the visibility of embedded circumstellar disks. The procedure for defining disks and proplyds boundaries assumed in this paper is systematic and precise but dependent on the surrounding object nebular intensity, which is highly variable. Nevertheless, this is still the best possible size measurement procedure. Previous papers do not define the measuring procedure and proplyds’ sizes are even more dependent on personal judgment. Estimating the completeness of our sample is a difficult task. Still, given that the spatial resolution of the WFPC2 image is 67.5 AU, we can crudely estimate that to a visual extinction limit of a few magnitudes of A<sub>V</sub> our sample should be essentially complete for disk diameters larger than $``$ 100–150 AU.
### 3.7 Size distribution of disks
In Figure 4, we present the spatial diameter distribution of all the 149 disks in our sample of the Trapezium cluster. The “star” symbols represent the 6 brightest Trapezium OB stars. The filled circles indicate the positions of the 135 bright proplyds’ disks; gray for the 10 proplyds with embedded silhouette disks and black for the remain 125 from which the disk diameter is 49% of the correspondent IF diameter. The unfilled circles represent the positions of the 14 pure silhouettes. The diameters of the circles are proportional to the disk diameters they represent. The pure silhouettes are the largest disks that lie at relatively large projected distances from the Trapezium. Two proplyds seem to fall off the size distribution i) 114-426, with a diameter of 1242 AU, is by far the largest and ii) 218-354, with a diameter of 675 AU. The dashed circle indicates a radial distance of 0.3 pc from $`\theta ^1`$ Ori C and marks the limiting border of the FUV-dominated region of the Störzer & Hollenbach storzer99 (1999) photoevaporation model. It is also the distance between the two ionizing stars, $`\theta ^1`$ Ori C and $`\theta ^2`$ Ori A, with the EUV photoionizing luminosity of the former 3-4 times greater than the later (O’Dell 2001b). About 73% of the proplyds in our sample are inside this region and 60% at a projected distance less than 0.2 pc. The proplyd situated at the largest distance from $`\theta ^1`$ Ori C is 005-514 at $``$ 264<sup>′′</sup> or 0.58 pc. We do not find a correlation between the disk diameters and the projected distance to $`\theta ^1`$ Ori C, as already discussed before.
Figure 5 represents the disk diameter distribution histogram for the total sample of 149 proplyds: 135 bright proplyds, in gray, and 14 silhouettes, in black. To select a statistically significant binning for the distribution we used a density kernel estimator (Silverman 1996) that indicated a bin-width of 100 AU at a 90% confidence limit. The resultant disk diameter distribution seems unimodal, with a tail to the large diameters, suggesting a single population of disks that is well characterized by a power-law distribution.
Figure 6 represents the power-law fit, in a log-log scale, to the sample of 149 disks presented in Figure 5. The first and incomplete bin (0-100 AU) was deliberately excluded from the power-law fit. We assumed that all the disk diameters (excluding silhouettes) are 49 $`\pm `$ 16% of the proplyds’ ionization front diameters, based on the 10 well-known cases with resolved embedded disks. The data is well described by a power-law with an exponent of $``$1.9 $`\pm `$ 0.3. To determine the uncertainty on the power-law fit we performed several tests, shifting the binning starting point by 25, 50 and 75 AU and the diameter ratio R for the 125 bright proplyds, to R = 33% and 65%, accounting for the dispersion found in this ratio. The power-law exponent is the average of the 12 coefficients determined in these tests. An obvious and unique deviation to the power-law distribution characterizing our sample is the silhouette 114-426. With a diameter of about 1242 AU this large disk seems to be an exception compared to other proplyd disks in Orion. About 80 to 90% of Trapezium’s young stars and $``$ 65% of the brown dwarfs show infrared excess emission characteristic of circumstellar disks (Lada et al. 2000; Muench et al. 2001). Assuming that the fraction of disks is not a strong function of depth into the Orion molecular cloud (near-infrared surveys used by Lada et al. lada00 (2000) and Muench et al. muench00 (2000) are less affected by dust extinction than the H<sub>α</sub>/HST image analyzed in this paper) and given that the HST images reveal, for about the same area of the cluster, that only $``$ 50% of the sources are associated with extended circumstellar structures, one can infer that a relatively large fraction of sources in the HST images (30 to 40%) also have disks but they are smaller than the resolution limit of the HST image. Since there were $``$ 300 YSO imaged with the HST then $``$ 85%, or about 255 of them, should have circumstellar disks but we only identify 149. This means that the other 106 disks are probably too small to have been resolved by the HST. The virtual sample, or corrected sample of disks in the Trapezium, is composed by the 149 disks from our sample plus 106 disks of 50 AU added. The dashed line in Fig. 5 represents the histogram for the virtual sample of 255 circumstellar disks or 85% of the YSO observed in the Trapezium cluster. The cumulative disk diameter function of the histogram in Fig. 5 is given by Fig. 7 and indicates that 75 to 80% of disks have diameters smaller than 150 AU (the uncorrected value is about 60%). This result is not consistent with Rodmann rodmann02 (2002) who finds that 90% of the disks have diameters smaller than 80 AU (as referenced in Bate et al. 2003). From this cumulative function we can estimate that about 40% of the disks in the Trapezium have radius larger than 50 AU, the Solar System size.
### 3.8 Does disk size depend on the parent star mass?
In Figure 8 we present disk diameter as a function of central star spectral-type. Our proplyd list was compared with the Lynne Hillenbrand’s online working version of table1 (last update: September 2003) from Hillenbrand hillenbrand97 (1997)<sup>4</sup><sup>4</sup>4http://www.astro.caltech.edu/$``$lah/papers/$`orion_{}`$main.table1.working. The selection criterion consisted in the overlap of the coordinates of the proplys in the two tables with radius less than 2 arcsec. For the resulting 52 sources in common, we considered the spectral-types derived by Hillenbrand hillenbrand97 (1997) (optical spectroscopy) and Luhman et al. luhman00 (2000) (NIR spectroscopy). Only 43 of the 52 sources had a spectral classification and only 21 of them are classified by Luhman et al. luhman00 (2000). For the 21 Luhman sources, represented in Figure 8 as gray squares, we assumed the error to be equal to the uncertainty interval in Luhman’s NIR spectral classification. The errors assumed for the 22 Hillenbrand sources are $`\pm `$1 spectral subclass, for spectral types later than K7, and $`\pm `$1/2 class for spectral types former than K7, as described in her paper. They are represented as black triangles in the plot. A systematic error of 45 AU or 1 pixel was considered in the disk diameters.
In order to describe the strength of the linear association between the mass of the star and the diameter of the disk surrounding it, several tests were performed on the data. We calculated the Pearson linear correlation coefficient for all the 43 sources of overlap between our sample and Hillenbrand’s catalog. To verify if this linear correlation was or not dependent on the way that the YSOs spectral types were determined, we separated the 43 combined data in two samples: 21 proplyds from Luhman et al. luhman00 (2000), and 22 proplyds from Hillenbrand hillenbrand97 (1997). For the combined sample we derive a coefficient of 0.3 indicating a small correlation between the two variables (10%). However, we get a better correlation for Luhman sources alone (suggesting a significant correlation between diameters and masses) than in Hillenbrand sample where we get a very small correlation. A larger sample or a better understanding of spectral typing differences between optical and NIR should clarify this inconsistency. We conclude that, over the sampled mass range (late G to late M stars) and with the present spectral data, there is not an obvious correlation between disk diameter and stellar mass.
## 4 Discussion
An important result of this work is that there seems to exist one single population of disks well characterized by a power-law distribution. Albeit the young age of the Trapezium, and given that disk destruction is well underway, it is perhaps too late to tell if the present day disk size distribution is primordial or if it is a consequence of the massive star formation environment. Nevertheless, the simple existence of a well defined disk size distribution hints at the existence of a primordial disk size distribution. Intuitively one would argue that a disk initial mass and size would be proportional to stellar mass (continuum surveys performed in the mm/submm domain suggest a tendency for the disk mass to increase with the stellar mass, e.g., Natta et al. 2000) but the results expressed in Figure 8, that there is not an obvious correlation between disk size and stellar mass, are intriguing and suggest a more complex picture. Most likely there is a random disk destruction process, or a combination of processes, that could be responsible for this lack of correlation.
The most important disks destruction processes are 1) viscous accretion, 2) close stellar encounters, and 3) stellar winds. Several environment related disk destruction mechanisms (2 and 3) have recently been proposed in the literature: Bate et al. bate03 (2003) performed sophisticated hydrodynamical simulations of cluster formation and find that most circumstellar disks are severely truncated by dynamical encounters. On the other hand, Richling & Yorke richling00 (2000), considering photoevaporation, high initial disk masses, and fixed distances to $`\theta ^1`$ Ori C, obtain the present radius and masses observed in Trapezium’s proplyds, for a timescale consistent with the age of $`\theta ^1`$ Ori C ($``$ 0.5 Myr). Also, Scally & Clarke scally01 (2001) performed numerical N-body simulations of the Orion Nebula cluster and determine a very low probability of an encounter at a d $`<`$ 100 AU, at the cluster’s present age, suggesting photoevaporation as the most significant disk destruction mechanism. If significant UV radiation is available, photoevaporation is the dominant disk dispersal mechanism (Hollenbach et al. 2000). The intense UV radiation field heats the disk surface, drives mass-loss and produces the bright ionization fronts, on a relatively short timescale. For example, the Trapezium proplyds have radial velocities of 24 to 30 Kms<sup>-1</sup> (Henney & O’Dell 1999) and inferred disk masses of 0.005-0.02 M (Lada et al. 1996; Bally et al. 1998b), with the 0.02 M upper bound determined for the largest pure silhouette 114-426. Mass-loss rates of order 10<sup>-7</sup> to 10<sup>-6</sup> Myr<sup>-1</sup> have been measured for a number of proplyds by several groups (Churchwell et al. 1987; Johnstone et al. 1997; Bally et al. 1998b; Johnstone et al. 1998; Henney & Arthur 1998; Henney & O’Dell 1999) with a mean mass-loss rate of 3.3 $`\times `$ 10<sup>-7</sup> Myr<sup>-1</sup> for all 31 proplyds that have been studied with sufficient detail (out of a total of $``$ 150 Orion bright proplyds). This implies that a minimum Solar Nebula mass disk (0.01 M) will loose half of its mass in 10<sup>4</sup> to 10<sup>5</sup> years and even less for the less massive observed objects. Such lifetimes are quite small compared to the estimated age of $`\theta ^1`$ Ori C or 0.5 Myr. Herbig & Terndrup herbig86 (1986) and Hillenbrand hillenbrand97 (1997) find that the ages for low-mass stars in the Orion Nebula cluster range from 3 x 10<sup>5</sup> to 10<sup>6</sup> yr, with only a few as young as 10<sup>5</sup> yr. This brings a major conundrum, since proplyds exist and have been surviving the photoevaporation. There are two possible explanations: either $`\theta ^1`$ Ori C is very young, or the illumination time of the proplyds is short (which would require a large spatial motion of $`\theta ^1`$ Ori C). Recently, Tan tan04 (2004) claimed that the BN object is a runaway B star dynamical ejection event by the $`\theta ^1`$ Ori C star, just $``$ 4000 years ago. Regardless, the point is that photoevaporation is likely to be the dominant disk destruction process. Still, relevant observational studies of the dynamics of embedded stellar clusters are not available yet (efficient near-infrared high-resolution spectrographs are only now becoming available) and it is too soon to dismiss close stellar encounters as an important disk destruction process in young clusters.
### 4.1 The Trapezium as the birthplace of Solar system analogues
Can planet formation endure disk destruction mechanisms?
If the growth of larger particles can occur before the removal of gas and small particles, planets may nevertheless form inside the dust disks embedded in such an adverse environment. There are evidences for grain growth in the largest silhouette 114-426 (Throop et al. 2001; Shuping et al. 2003). Recently, Throop & Bally throop05 (2005) proposed that UV radiation can stimulate the rapid formation of planetesimals in externally-illuminated protoplanetary disks. It might then be that photoionization has an overall positive feedback on planet formation, while being an important (gas) disk destruction process. If this is the case, this would have far reaching implications to planet formation in the Galaxy, and likely in the Universe, since most stars are born in clusters containing O and B stars (Lada & Lada 2003). In the specific case of the Trapezium, where about half of the disks are larger or about the Solar system size (see Section 3.7), one could expect planetary systems to be common and not fundamentally different from the Solar system itself.
It is striking that the remnant Solar system disk appears to be truncated, not unlike the Trapezium disk discussed in this paper. It is now a rather well established and intriguing fact that the Kuiper belt, which is a repository of the solar system’s most primitive matter, has a well defined outer edge at about 50 AU (e.g., Jewitt 2002). At least three mechanisms for its origin have been proposed, none of which has raised the general consensus of the community of the experts (see Morbidelli et al. 2003 for a review): 1) a dynamical origin involving a 1 Gyr lived Mars mass at object at 60 AU (Brunini & Melita 2002), 2) differential accretion rates producing a disk edge (Weidenschilling 2003), and 3) planetesimal disk truncation by the passage of a star in the vicinity of the Sun (e.g., Ida et al. 2000, Kobayashi & Ida 2001). The HST images of the Trapezium cluster disks, in particular the silhouette disks, present us with clear evidence that these disks seem to have already well defined edges at their relatively very young ages. If the Sun was born in a stellar cluster, the probable case as most stars are born in clusters, this is suggestive that the origin of the Kuiper belt outer edge is likely to be due to the star formation environment and disk destruction processes (photoevaporation, collisions) present in archetypal star formation factories such as the Trapezium. If this is the case, these well defined outer edges are imprinted in these disks at the earliest phases of their evolution, even before the formation of Kuiper belt-like objects.
The results in this paper, together with future higher-resolution observations and modeling, can help identify the dominant process for circumstellar disk destruction and provide insights into the survival rate of circumstellar material surrounding the YSOs, and therefore, insights on star and planet formation in very young clusters containing O and B stars, the typical nursery for most stars in the Galaxy.
## 5 Conclusions
The main results in this paper can be summarized as follow:
* There is no meaningful correlation between the IF “chord diameters” and the projected distance to $`\theta ^1`$ Ori C. We find the same tendency in our analysis of previously published data.
* Direct measurements of 10 disks embedded in bright proplyds show a great dispersion in the diameter ratio R = disk diameter/ IF diameter, and R seems not to be correlated with the projected distance to $`\theta ^1`$ Ori C. 244-440 (R = 15%) and 072-135 (R = 71%) illustrate the 2 opposite “extreme” cases.
* Assuming R = 49 $`\pm `$ 16% (the average ratio from the 10 cases) to compute the disk sizes for the 125 bright proplyds we determined an unimodal disk size distribution, representing a single population of disks, that is well characterized by a power-law distribution with exponent of -1.9 $`\pm `$ 0.3.
* For the stellar mass sampled (from late G to late M stars) we find that there is no obvious correlation between disk size and stellar mass.
* The pure silhouettes are clearly the largest disks and have large projected distances to $`\theta ^1`$ Ori C. In particular, object 114-426 is rather unique given its size and it falls off the above characterized disk size distribution.
* We estimate that about 40% to 45% of the Trapezium cluster disks have radius larger than 50 AU, the Solar System size.
* We suggest that the origin of the Solar system’s (Kuiper belt) outer edge is likely to be due to the star formation environment and disk destruction processes (photoevaporation, collisions) present in the stellar cluster on which the Sun was probably formed.
* We identified a previously unknown proplyd and named it 266-557, following convention.
This statistical analysis should be repeated in function of ”true” distances to the OB stars and not the projected ones. Proplyd geometry and true distances will lead to the 3-D spatial distribution of proplyds in Orion.
###### Acknowledgements.
The authors are grateful to Nicole Homeier, Michael Liu, Martino Romaniello, Herve Bouy and Ricardo Demarco for fruitful discussions.
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# Property (𝑇) and rigidity for actions on Banach spaces
## 1. Introduction and the Main Results
### 1.a.
Since its introduction by Kazhdan in \[Ka\], property $`(T)`$ became a fundamental concept in mathematics with a wide range of applications to such areas as:
• The structure of infinite groups – finite generation and finite abelanization of higher rank lattices \[Ka\], obstruction to free or amalgamated splittings \[Wa\], \[A\], \[M4\] structure of normal subgroups \[M2\] etc.;
• Combinatorics – the first construction of expanders \[M1\] (see \[Lu\]);
• Operator algebras – factors of type II<sub>1</sub> whose fundamental group is countable \[C\] or even trivial \[Po1\]); rigidity theorems for the factors associated to Kazhdan group \[Po2\];
• Ergodic theory – rigidity results related to Orbit Equivalence \[Po3\],\[Hj\]; the Banach–Ruziewicz problem \[M3\],\[Su\];
• Smooth dynamics – local rigidity \[FM1\],\[FM2\]; actions on the circle \[N1\] (and \[PS\],\[Rz\]).
It has also been an important tool in providing interesting (counter) examples: to Day’s “von Neumann conjecture” \[Gr1, 5.6\] and in the context of the Baum–Connes conjecture \[HLS\] (related to \[Gr2\]).
Initially defined in terms of unitary representations, property $`(T)`$ turned out to be equivalent to Serre’s property $`(FH)`$ – a fixed point property for affine isometric actions on Hilbert spaces that can be rephrazed as cohomological vanishing. (The equivalence holds for $`\sigma `$-compact groups, in particular all locally compact second countable groups, and was proved by Delorme \[D\] and Guichardet \[Gu\]. As pointed out by Y. de Cornulier, uncountable discrete groups that have G. Bergman’s cofinality property \[Bn\] have $`(FH)`$ but fail $`(T)`$.) Some of the above applications use this latter characterization. Recently Shalom \[Sh\] described the *reduced* 1-cohomology with unitary coefficients for irreducible lattices in products of completely general locally compact groups. This led to a list of new rigidity results and added such lattices to the list of “naturally rigid” groups. For further details and more references on these topics, we suggest the monography \[HV\] and the forthcoming \[BHV\].
### 1.b.
Motivated by these broad themes: property $`(T)`$, property $`(FH)`$, lattices in semisimple groups and in general products, we study similar notions in the broader framework of Banach spaces rather than Hilbert spaces. Some of the results below apply to general *superreflexive* Banach spaces, whilst some are specific to the subclass of $`L^p(\mu )`$-spaces with $`1<p<\mathrm{}`$. (A Banach space is superreflexive if it admits an equivalent uniformly convex norm, see Proposition 2.3 below.)
One of the motivations to consider such questions came from the work of Fisher and Margulis \[FM1\], \[FM2\], in which an $`L^p`$ analogue of property $`(T)`$ with $`p2`$ allowed them to weaken smoothness assumptions in their results.
The harder question of fixed point results for *affine actions* on $`L^p`$ for $`p2`$ (see Theorem B below) has applications e.g. for actions on the circle \[N2\], \[BHV\].
### 1.c.
Let $`G`$ be a topological group and $`B`$ a Banach space. By a *linear isometric $`G`$-representation* on $`B`$, we shall mean a continuous homomorphism $`\varrho :G𝐎(B)`$ where $`𝐎(B)`$ denotes the (“orthogonal”) group of all invertible linear isometries $`BB`$ (see Lemma 2.4 for a clarification of the continuity assumption). We say that such a representation *almost has invariant vectors* if
(1.i)
$$\text{compact subset }KG,\underset{v=1}{inf}\mathrm{diam}(\varrho (K)v)=0.$$
Denote by $`B^{\varrho (G)}`$ the closed subspace of $`G`$-fixed vectors; the $`G`$-representation $`\varrho `$ descends to a linear isometric $`G`$-representation $`\varrho ^{}`$ on $`B^{}=B/B^{\varrho (G)}`$ (see Remark 2.11 for more details in the case of superreflexive spaces). We shall use the following as a Banach space analogue of Kazhdan’s property $`(T)`$:
###### Definition 1.1.
Let $`B`$ be a Banach space. A topological group $`G`$ is said to have property $`(T_B)`$ if for any continuous linear isometric $`G`$-representation $`\varrho :G𝐎(B)`$ the quotient $`G`$-representation $`\varrho ^{}:G𝐎(B/B^{\varrho (G)})`$ does not almost have $`G`$-invariant vectors.
Note that if $`B`$ is a Hilbert space, $`\varrho ^{}`$ is isomorphic to the restriction of $`\varrho `$ to the orthogonal complement $`(B^{\varrho (G)})^{}`$ of the subspace of $`\varrho (G)`$-invariants. Thus for Hilbert spaces the above definition agrees with Kazhdan’s property $`(T)`$.
Let $`\mu `$ be a $`\sigma `$-finite measure on a standard Borel space $`(X,)`$. We are most interested in the family $`L^p(\mu )`$, $`1<p<\mathrm{}`$, of Banach spaces, which are close relatives of Hilbert spaces. They also possess a rich group of linear isometries $`𝐎(L^p(\mu ))`$.
###### Theorem A.
Let $`G`$ be a locally compact second countable group. If $`G`$ has Kazhdan’s property $`(T)`$ then $`G`$ has property $`(T_B)`$ for Banach spaces $`B`$ of the following types:
* $`L^p(\mu )`$ for any $`\sigma `$-finite measure $`\mu `$ and any $`1p<\mathrm{}`$.
* A closed subspace of $`L^p(\mu )`$ for any $`1<p<\mathrm{}`$, $`p4,6,8,\mathrm{}`$.
* A quotient space of $`L^p(\mu )`$ for any $`1<p<\mathrm{}`$, $`p\frac{4}{3},\frac{6}{5},\frac{8}{7},\mathrm{}`$.
If $`G`$ has $`(T_{L^p([0,1])})`$ for some $`1<p<\mathrm{}`$ then $`G`$ has Kazhdan’s property $`(T)`$.
### 1.d.
Next we consider group actions by isometries on Banach spaces. By the Mazur–Ulam theorem, such actions are always affine with the linear part being isometric as well (working with *real* Banach spaces).
###### Definition 1.2.
We say that $`G`$ has property $`(F_B)`$ if any continuous action of $`G`$ on $`B`$ by affine isometries has a $`G`$-fixed point.
When $`B`$ is a Hilbert space this is precisely Serre’s property $`(FH)`$. Delorme \[D\] and Guichardet \[Gu\] proved that properties (T) and $`(FH)`$ are equivalent for $`\sigma `$-compact groups. Below we summarize the relations between properties $`(T)`$ and $`(F_B)`$ which hold for general groups.
###### Theorem 1.3.
For a locally compact second countable group $`G`$ we have
1. $`(F_B)`$ implies $`(T_B)`$ for any Banach space $`B`$.
2. $`(T)`$ implies $`(F_B)`$ for closed subspaces $`B`$ of $`L^p(\mu )`$ where $`1<p2`$.
Likewise for subspaces of $`L^1`$ and of the pseudo-normed spaces $`L^p(\mu )`$,
$`0<p<1`$, except one obtains only bounded orbits instead of fixe pointsSee Example 2.23 for an example without fixed point..
3. $`(T)`$ also implies $`(F_B)`$ for closed subspaces of $`L^p(\mu )`$ for $`2p<2+\epsilon `$, where $`\epsilon =\epsilon (G)>0`$ might depend on the Kazhdan group $`G`$.
###### Remarks 1.4.
(1) is essentially due to Guichardet \[Gu\] as his proof of $`(FH)(T)`$ applies to all Banach spaces. We give two proofs for (2) reducing the problem, in both, to one of the proofs of $`(T)(FH)`$. We note that the particular case of $`p=1`$ in (2) is one of the results of \[RS\]. Statement (3) is due to Fisher and Margulis (unpublished). With their kind permission we have included their argument here (see Section 3.c).
The above results imply that any locally compact group $`G`$ with Kazhdan’s property $`(T)`$ has property $`(T_{L^p})`$ for all $`1<p<\mathrm{}`$, and has the fixed point property $`(F_{L^p})`$ for $`1<p<2+\epsilon (G)`$. It turns out, however, that many Kazhdan groups (e.g. hyperbolic ones) do not have property $`(F_{L^p})`$ for large values of $`p`$.
Indeed, in his study of $`L^p`$-cohomology, Pansu \[Pa\] proved that $`\mathrm{𝐒𝐩}_{n,1}(𝐑)`$ and cocompact lattices in these groups have a non-trivial first $`L^p`$-cohomology $`L^pH^1`$ for all $`p>4n+2`$. This is equivalent to asserting that for $`p>4n+2`$ these groups admit fixed-point-free affine isometric actions on $`L^p(G)`$ with linear part being the regular representation. Hence these groups do not have property $`(F_{L^p})`$ for $`p>4n+2`$, whilst enjoying $`(T)`$.
More generally, $`L^pH^1(\mathrm{\Gamma })`$ and hence $`H^1(\mathrm{\Gamma },\mathrm{}^p\mathrm{\Gamma })`$ is non-zero for any non-elementary hyperbolic group when $`p`$ is large enough. Indeed, Bourdon and Pajot identify this cohomology with a Besov space of functions on the boundary, which they prove to be non-trivial as soon as $`p`$ is larger than the Hausdorff dimension of an Ahlfors-regular metric on the boundary, see Corollaire 6.2 in \[BP\]. Again, this contradicts $`(F_{L^p})`$ for large $`p`$.
More recently, using Mineyev’s homological bicombings \[Mi\], Yu \[Y\] gave a very short proof that any hyperbolic group $`\mathrm{\Gamma }`$ admits a *proper* action by affine isometries on $`\mathrm{}^p(\mathrm{\Gamma }\times \mathrm{\Gamma })`$ if $`p`$ is large enough. This is a strong negation of $`(F_{L^p})`$ for hyperbolic groups and all their infinite subgroups. The corresponding strenghtening of the above mentioned \[Pa\], \[BP\] for rank one Lie (or algebraic) groups $`G`$ has been established by Cornulier–Tessera–Valette in \[CTV\]: For any $`p>1`$ larger than the Hausdorff dimension of the boundary, there is a *proper* affine isometric action on $`L^p(G)`$ whose linear part is the regular representation. In particular, this holds for $`\mathrm{𝐒𝐩}_{n,1}(𝐑)`$ when $`p>4n+2`$.
### 1.e.
Our next goal is now, by contrast, to establish $`(F_{L^p})`$ for certain groups. It is often remarked that property $`(T)`$ for (simple) higher rank Lie groups and their lattices is more robust than property $`(T)`$ enjoyed by the rank one groups $`\mathrm{𝐒𝐩}_{n,1}(𝐑)`$ and many other Gromov hyperbolic groups. In view of the preceding discussion of hyperbolic groups and $`\mathrm{𝐒𝐩}_{n,1}(𝐑)`$, the following result might be viewed as yet another evidence supporting this view.
###### Theorem B.
Let $`G=_{i=1}^m𝐆_i(k_i)`$, where $`k_i`$ are local fields (of any characteristic), $`𝐆_i(k_i)`$ are $`k_i`$-points of Zariski connected simple $`k_i`$-algebraic groups $`𝐆_i`$. Assume that each simple factor $`𝐆_i(k_i)`$ has $`k_i`$-rank $`2`$.
Then $`G`$ and the lattices in $`G`$ have property $`(F_B)`$ for all $`L^p(\mu )`$-related spaces $`B`$ as in (i)–(iii) in Theorem A, assuming $`1<p<\mathrm{}`$.
### 1.f.
A broader class of spaces in which we propose to study properties $`(T_B)`$ and $`(F_B)`$ consists of *superreflexive* spaces, which can be defined as topological vector spaces isomorphic to uniformly convex Banach spacesFor spaces that are only *strictly convex*, the fixed point property always fails \[BG\],\[HP\].. In this context we consider linear representations (resp. affine actions) which are *uniformly equicontinuous*; more concretely, for any given norm compatible with the topology, the class of all such linear representations (resp. affine actions) is that of *uniformly bounded* linear representations (resp. *uniformly Lipschitz* affine actions). It turns out that such representations (resp. actions) can always be viewed as isometric with respect to some equivalent norm that is *simultaneously* uniformly convex and uniformly smooth (Proposition 2.13).
Note that whether a given linear $`G`$-representation almost contains invariant vectors or not, in the sense of (1.i), does not depend on a particular norm among all mutually equivalent norms. Hence we can make the following
###### Definition 1.5.
Let $`B`$ be a superreflexive topological vector space and $`G`$ a locally compact second countable group. We say that $`G`$ has property $`(\overline{T}_B)`$ if for every uniformly equicontinuous linear representation $`\varrho `$ of $`G`$ on $`B`$ the quotient $`G`$-representation on $`B/B^{\varrho (G)}`$ does not almost have invariant vectors.
Likewise, $`G`$ has $`(\overline{F}_B)`$ if every uniformly equicontinuous affine $`G`$-action on $`B`$ has a fixed point.
###### Conjecture 1.6.
Higher rank groups $`G=𝐆_i(k_i)`$ as in Theorem B and their lattices have property $`(\overline{F}_B)`$, and hence $`(\overline{T}_B)`$, for all superreflexive $`B`$.
###### Remark 1.7.
To support this conjecture let us point out the following:
(1) Much of our proof of Theorem B is done in the broad context of uniformly equicontinuous affine actions on general superreflexive spaces except for one argument – a version of relative property $`(T_B)`$, whose proof is special to $`L^p`$-related spaces.
(2) V. Lafforgue proved \[Lg\] that the group $`\mathrm{𝐏𝐆𝐋}_3(𝐐_p)`$ has property $`(\overline{T}_B)`$ for all superreflexive $`B`$ (his result is actually stronger, in that he allows linear representations with *slowly growing*, rather *uniformly bounded* Lipschitz norms, see Theorem 3.2, Definition 0.2 and the discussion preceding it in \[Lg\]). Combined with our proof of Theorem B it implies for example that $`\mathrm{𝐒𝐋}_n(𝐐_p)`$, $`n4`$, has property $`(\overline{F}_B)`$.
(3) Y. Shalom has proved (unpublished) that for Hilbert spaces $``$ higher rank groups (and their lattices) have property $`(\overline{F}_{})`$, and hence $`(\overline{T}_{})`$, whilst rank one groups have neither $`(\overline{F}_{})`$ nor $`(\overline{T}_{})`$.
### 1.g.
One way to generalize the context of semisimple (non-simple) Lie/algebraic groups is simply to consider general products $`G=G_1\times \mathrm{}\times G_n`$ of $`n2`$ *arbitrary* topological groups. In the absence of any assumption on the factors $`G_i`$, one can still establish splitting results for uniformly equicontinuous affine $`G`$-actions on superreflexive spaces.
###### Theorem C.
Let $`G=G_1\times \mathrm{}\times G_n`$ be a product of topological groups with a continuous action by uniformly equicontinuous affine maps on a superreflexive topological vector space $`B`$ without $`G`$-fixed point. Assume that the associated linear $`G`$-representation $`\varrho `$ does not almost have non-zero invariant vectors.
Then there is a $`G`$-invariant closed complemented affine subspace $`\underset{¯}{B}B`$ and an affine equicontinuous $`G`$-equivariant isomorphism $`\underset{¯}{B}B_1\mathrm{}B_n`$, where each $`B_i`$ is a superreflexive Banach space with an equicontinuous affine $`G`$-action factoring through $`GG_i`$.
###### Remarks 1.8.
(1) If $`G`$ has property $`(\overline{T}_B)`$ then the assumption that $`\varrho `$ does not almost have invariant vectors is redundant.
(2) In the particular case where $`B`$ is a Hilbert space and $`G`$ locally compact acting by affine isometries, a stronger result was established by Shalom in \[Sh\]: One assumes only that the *affine* $`G`$-action does not almost have fixed points. We replace Shalom’s Hilbertian approach with an analogue of the geometric method used in the splitting theorem of \[Mo2\].
(3) This result can be reformulated in terms of the cohomology of the associated linear $`G`$-representation $`\varrho `$ on $`B`$ as
$$H^1(G,B)\underset{i=1}{\overset{n}{}}H^1(G_i,B^{\varrho (_{ji}G_j)}).$$
It should be stressed that no such product formula holds in general. Not only does it fail for more general Banach spaces (Example 2.27), but even for Hilbert space one needs at least Shalom’s assumption mentioned above. Compare the similar situation for the cohomological product formulas of \[Sh\] and \[BMd\].
### 1.h.
When $`G`$ is locally compact, we can as in the Lie case consider its lattices. One then calls a lattice $`\mathrm{\Gamma }<G`$ *irreducible* if its projections to all $`G_i`$ are dense. The above Theorem C can be used to establish a superrigidity result for irreducible lattices much in the way of \[Sh\]. (The general idea to use irreducibility in order to transfer results from $`G_1\times \mathrm{}\times G_n`$ to $`\mathrm{\Gamma }`$ was also illustrated in \[BMz\], \[BMd\],\[MS\]; it seems to originate from the work of Margulis and \[BK\]; lattices in products of completely general locally compact groups were first studied by Shalom \[Sh\].)
###### Theorem D.
Let $`\mathrm{\Gamma }`$ be an irreducible uniform lattice in a locally compact $`\sigma `$-compact group $`G=G_1\times \mathrm{}\times G_n`$. Let $`B`$ be a superreflexive space with uniformly equicontinuous affine $`\mathrm{\Gamma }`$-action. Assume that the associated linear $`\mathrm{\Gamma }`$-representation does not almost have invariant vectors.
Then there is a $`\mathrm{\Gamma }`$-closed complemented affine subspace of $`B`$ on which the $`\mathrm{\Gamma }`$-action is a sum of actions extending continuously to $`G`$ and factoring through $`GG_i`$.
###### Remark 1.9.
More precisely, the conclusion means that there are superreflexive spaces $`E_i`$ endowed each with a continuous uniformly equicontinuous affine $`G`$-action factoring through $`GG_i`$ and a $`\mathrm{\Gamma }`$-equivariant affine continuous map $`_{i=1}^nE_iB`$. Equivalently, the cocycle $`b:\mathrm{\Gamma }B`$ of the original $`\mathrm{\Gamma }`$-action is cohomologous to a sum $`b_1+\mathrm{}+b_n`$ of cocycles $`b_i`$ ranging in a subspace $`B_iB`$ on which the *linear* $`\mathrm{\Gamma }`$-representation extends continuously to a $`G`$-representation factoring through $`G_i`$ and such that $`b_i`$ extends continuously to a cocycle $`GG_iB_i`$ (with respect to the corresponding linear $`G`$-representation). Moreover, $`B_iE_i`$ as $`G`$-spaces.
If one disregards a component of $`B`$ where the linear $`\mathrm{\Gamma }`$-representation ranges in a compact group of operators, this *sum of actions* is actually just a direct sum $`B_iB`$ (see Remark 8.10).
###### Remark 1.10.
A *uniform lattice* (in a locally compact group) is just a discrete cocompact subgroup; the theorem however also holds for certain non-uniform lattices, see Section 8 (Theorem 8.3). Similar arguments allow us to generalise slightly Shalom’s *superrigidity for characters*, see Theorem 8.4.
### Organization of the Paper
In Section 2 we collect preliminary facts and lemmas on uniformly convex/smooth and superreflexive Banach spaces, linear representations and affine isometric on such spaces, special properties of $`L^p`$-spaces, and some general remarks and basic counter-examples. In Section 3 Theorem 1.3 is proved. Equivalence of properties $`(T)`$ and $`(T_{L^p})`$ (Theorem A) is proved in Section 4. In Section 5 we discuss higher rank groups and prove Theorem B. Section 6 studies minimal convex sets. Section 7 addresses product groups and proves the splitting theorem (Theorem C); it also proposes a proof of Theorem B that provides some evidence for Conjecture 1.6. In Section 8, we prove Theorem D. Appendix 9 describes Shalom’s proof of a generalized Howe–Moore theorem.
### Acknowledgments
We would like to thank D. Fisher and G.A. Margulis for their interest in this work and for letting us include their argument for $`2<p<2+\epsilon `$ in Section 3.c. We are indepted to A. Naor for several helpful conversations and to A. Nevo for his remarks on the first manuscript. We are grateful to Y. Shalom for letting us give his proof of a ucus version of Howe–Moore (Theorem 9.1).
## 2. Preliminaries
This section contains basic definitions, background facts and some preliminary lemmas to be used in the proofs of our main results.
### 2.a. Banach Spaces
Let $`B`$ be a Banach space; unless otherwise specified, we take the reals as scalar field. We denote by $`\mathrm{S}(B)=\{xB:x=1\}`$ its unit sphere. For $`xB`$ and $`r>0`$ we denote by $`\mathrm{B}(x,r)`$ and $`\overline{\mathrm{B}}(x,r)`$ the open, respectively closed, ball of radius $`r`$ around $`x`$.
A Banach space $`B`$ is said to be *strictly convex* if its unit sphere does not contain straight segments, or equivalently if $`(x+y)/2<1`$ whenever $`xy\mathrm{S}(B)`$. A Banach space $`B`$ is called *uniformly convex* if the *convexity modulus* function
(2.i)
$$\delta (\epsilon )=inf\{1x+y/2:x,y1,xy\epsilon \}$$
is positive $`\delta (\epsilon )>0`$ whenever $`\epsilon >0`$.
We shall also use the notion of *uniform smoothness* of Banach spaces, which is easiest to define as the uniform convexity of the dual space $`B^{}`$ (see \[BL, App. A\]). Hence a Banach space $`B`$ is *uniformly convex and uniformly smooth* (hereafter abbreviated *ucus*) if both $`B`$ and its dual $`B^{}`$ are uniformly convex.
###### Facts 2.1.
We refer to \[BL\] for the following:
1. The function $`\delta (\epsilon )`$ is non-decreasing and tends to $`0`$ when $`\epsilon `$ tends to $`0`$. If $`B`$ is uniformly convex then $`\delta (\epsilon )0\epsilon 0`$.
2. Uniformly convex Banach spaces are reflexive. Hence the class of ucus Banach spaces is closed under taking duals. This class is also closed under the operations of taking closed subspaces and quotients.
3. If $`B^{}`$ is strictly convex, in particular if $`B`$ is uniformly smooth, then every $`x\mathrm{S}(B)`$ has a unique supporting functional $`x^{}\mathrm{S}(B^{})`$, i.e. a unit functional with $`x,x^{}=1`$.
4. If $`B`$ is ucus then the *duality map* $`:\mathrm{S}(B)\mathrm{S}(B^{})`$, $`xx^{}`$, is a *uniformly continuous* homeomorphism with a uniformly continuous inverse.
5. To any non empty bounded subset $`EB`$ of a reflexive strictly convex Banach space $`B`$, one can associate a unique point $`C(E)B`$, the *circumcentre* of $`E`$ (a.k.a. the Chebyshev centre), defined as the unique $`xB`$ minimizing $`inf\{r>0:E\overline{\mathrm{B}}(x,r)\}`$.
The existence of $`x=C(E)`$ in (5) follows from weak compactness of closed bounded convex sets (i.e. from reflexivity), whilst the uniqueness follows from uniform convexity. Note that somewhat contrary to the intuition, it was shown by V. Klee \[Kl\] that if $`dim(B)3`$ and $`B`$ is not a Hilbert space, then there exist a bounded subset $`EB`$ for which $`C(E)`$ does not belong to the closed convex hull of $`E`$. The notion of circumcentre is also used in $`\mathrm{CAT}(0)`$ geometry. For CAT$`(0)`$ spaces, the circumcentre $`C(E)`$ always lies in the closed convex hull of $`E`$<sup>§</sup><sup>§</sup>§Note that Hilbert spaces are, in a sense, the most convex Banach spaces – they have the largest possible modulus of continuity $`\delta (\epsilon )`$ among Banach spaces. On the other hand, Hilbert spaces have the smallest possible modulus of continuity among CAT$`(0)`$ spaces. Thus, in a sense, CAT$`(0)`$ spaces are more convex then (non-Hilbertian) Banach spaces.
The following can be found e.g. in \[BL, A.6, A.8\]:
###### Theorem 2.2.
The following conditions on a topological vector space $`V`$ are equivalent:
1. $`V`$ is isomorphic to a uniformly convex Banach space.
2. $`V`$ is isomorphic to a uniformly smooth Banach space.
3. $`V`$ is isomorphic to a ucus Banach space.
The space $`V`$ is called *superreflexive* if these equivalent condition hold. The class of superreflexive spaces is closed under taking duals, closed subspaces and quotients of topological vector spaces.
### 2.b. Linear Representations
Let $`V`$ be a topological vector space. We denote by $`\mathrm{𝐆𝐋}(V)`$ the group of invertible linear transformations of $`V`$ which are continuous together with their inverses.
Following the standard terminology \[B1, Def. 2 of §2 no1\], a group $`G`$ of transformations of $`V`$ is *uniformly equicontinuous* (with respect to the uniform structure deduced from the topological vector space structure) if for any neighbourhood $`U`$ of $`0V`$ there exists a neighbourhood $`W`$ of $`0`$ such that
(2.ii)
$$xyWgG:g(x)g(y)U.$$
This definition will be applied to both linear groups, or affine groups.
For a topological vector space $`V`$, we denote by $`N(V)`$ the (a priori possibly empty) set of norms on $`V`$ defining the given topology. Elements of $`N(V)`$ will be called *compatible* norms and are pairwise equivalent.
The following key proposition is an equivariant version of Theorem 2.2. It enables us to reduce questions about uniformly equicontinuous linear representations on superreflexive spaces to isometric linear representations on ucus Banach spaces.
###### Proposition 2.3 (Invariant ucus norm).
For a superreflexive topological vector space $`V`$ and a group of linear transformations $`G`$ of $`V`$, the following conditions are equivalent:
1. $`G`$ is a uniformly equicontinuous group of linear transformations of $`V`$.
2. $`G`$ acts by uniformly bounded linear transformations with respect to any/all compatible norm on $`V`$.
3. $`G`$ acts by linear isometries with respect to some uniformly convex compatible norm on $`V`$.
4. $`G`$ acts by linear isometries with respect to some uniformly smooth compatible norm on $`V`$.
5. $`G`$ acts by linear isometries with respect to some uniformly convex and uniformly smooth compatible norm on $`V`$.
###### Proof.
The main part of the proof is the implication “\[(3) and (4)\]$``$(5)”; we begin by establishing this.
Let $`N(V)`$ denote the set of all compatible norms on $`V`$ equipped with the metric
$$d(_1,_2)=\underset{x0}{sup}\left|\mathrm{log}\frac{x_1}{x_2}\right|.$$
This is a complete metric space. Let $`N(V)^G`$ stand for the closed subspace of $`G`$-invariant norms in $`N(V)`$. Denoting by $`\delta _{}`$ the convexity modulus of $`N(V)^G`$, the subset $`N_{uc}(V)^G`$ of uniformly convex $`G`$-invariant norms on $`V`$ is given by the countable intersection
$$N_{uc}(V)^G=\underset{n=1}{\overset{\mathrm{}}{}}O_n,\text{where}O_n=\{N(V)^G:\delta _{}(1/n)>0\}.$$
Observe that the sets $`O_n`$ are open. If $`_0`$ is some fixed $`G`$-invariant compatible uniformly convex norm (given in (3)) then any $`N(V)^G`$ can be viewed as a limit of uniformly convex norms $`+\epsilon _0`$ as $`\epsilon 0`$. Hence $`N_{uc}(V)^G`$ is a dense $`G_\delta `$ set in $`N(V)^G`$.
By duality between $`N_{uc}(V^{})^G`$ and the set $`N_{us}(V)^G`$ of uniformly smooth norms in $`N(V)^G`$, the latter is also a dense $`G_\delta `$ set in the Baire space $`N(V)^G`$. In particular $`N_{uc}(V)^GN_{us}(V)^G`$ is not empty, as claimed.
Now we observe that “(1)$``$(2)” follows from the definitions and that “(5)$``$\[(3) and (4)\]” as well as “\[(3) or (4) or (5)\]$``$(2)” are trivial. Moreover, proving “(2)$``$(3)” will also yield “(2)$``$(4)” by duality, using the fact that the dual to a superreflexive space is superreflexive. Therefore it remains only to justify “(2)$``$(3)”:
Let $``$ be a compatible uniformly convex norm on $`V`$. The corresponding operator norms $`g=sup_{x0}gx/x`$ are uniformly bounded by some $`C<\mathrm{}`$. Hence
$$x^{}=\underset{gG}{sup}gx$$
defines a norm, equivalent to $``$, and $`G`$-invariant. It is also uniformly convex. Indeed, if $`x^{}=y^{}=1`$ and $`(x+y)/2^{}>1\alpha `$ then for some $`gG`$
$$(gx+gy)/2>1\alpha \text{whilst}gxx^{}=1,gyy^{}=1.$$
Thus $`\alpha \delta _{}(gxgy)\delta _{}(xy^{}/C)`$. Hence the convexity moduli satisfy
$$\delta _{^{}}(\epsilon )\delta _{}(\epsilon /C)>0\text{for all}\epsilon >0.$$
If $`G`$ is a topological group, one should impose a continuity assumption on linear $`G`$-representations on $`V`$, that is on homomorphisms $`\varrho :G\mathrm{𝐆𝐋}(V)`$. $`\mathrm{𝐆𝐋}(V)`$ is naturally equipped with the operator norm (which is too strong for representation theory), and with the *weak* and the *strong* operator topologies. For uniformly equicontinuous representations the latter two topologies impose the same continuity assumption:
###### Lemma 2.4.
Let $`G`$ be a topological group, $`V`$ a superreflexive topological vector space, and $`\varrho :G\mathrm{𝐆𝐋}(V)`$ a homomorphism. Then the following are equivalent.
1. $`\varrho `$ is weakly continuous.
2. $`\varrho `$ is strongly continuous.
3. The orbit maps $`g\varrho (g)u`$ are continuous.
4. The action map $`G\times BB`$ is jointly continuous.
Since there is an invariant complete norm on $`V`$, this is a special case of a well-known fact holding for all Banach spaces, see \[Mo1, 3.3.4\] for references. We give an elementary proof in the present case.
###### Proof.
Clearly it is enough to prove $`(1)(4)`$. Let $``$ be a $`\varrho (G)`$-invariant ucus norm on $`V`$. Assume $`g_neG`$ and $`u_nu\mathrm{S}(B)`$. Then
$`|\varrho (g_n)u_n,u^{}1|`$ $``$ $`\left|\varrho (g_n)u_n,u^{}\varrho (g_n)u,u^{}\right|+\left|\varrho (g_n)u,u^{}1\right|`$
$``$ $`u_nu+|\varrho (g_n)u,u^{}1|0`$
It follows that $`\varrho (g_n)u_nu`$ because
$$\frac{\varrho (g_n)u+u}{2},u^{}\frac{\varrho (g_n)u+u}{2}1\delta (\varrho (g_n)uu)$$
and the left hand side tends to $`1`$. ∎
### 2.c. Invariant complements
One of the convenient properties of Hilbert spaces is the existence of a canonical complement $`M^{}`$ to any closed subspace $`M`$. Recall that a closed subspace $`X`$ of a Banach space $`B`$ is called complemented if there is another closed subspace $`YB`$ such that $`B=XY`$ algebraically and topologically. This is equivalent to each of the following:
* There is a continuous linear projection from $`B`$ to $`X`$.
* There is a closed subspace $`Y`$ and a continuous linear projection $`p:BY`$ with $`\mathrm{ker}(p)=X`$.
A classical theorem of Lindenstrauss and Tzafriri says that every infinite dimensional Banach space which is not isomorphic to a Hilbert space, admits a non-complemented closed subspace \[LT2\]. However, for any uniformly equicontinuous linear representation $`\varrho `$ of a group $`G`$ on a superreflexive space $`B`$, the subspace of invariant vectors $`B^{\varrho (G)}`$ admits a canonical complement, described below.
In view of Proposition 2.3 we may assume that the representation is linear isometric with respect to a ucus norm on $`B`$, which allows to use the duality map of the unit spheres $`:\mathrm{S}(B)\mathrm{S}(B^{})`$.
Given any linear representation $`\varrho :G\mathrm{𝐆𝐋}(V)`$ there is an associated dual (or contragradient) linear $`G`$-representation $`\varrho ^{}:G\mathrm{𝐆𝐋}(V^{})`$ defined by
$$x,\varrho ^{}(g)y=\varrho (g^1)x,y(gG,xV,yV^{}).$$
If $`B`$ is a Banach space and $`\varrho :G𝐎(B)`$ is a linear isometric representation, then so is its dual $`\varrho ^{}:G𝐎(B^{})`$, where $`B^{}`$ is equipped with the dual norm. Hence the dual to a uniformly equicontinuous representation on a superreflexive space is also of the same type.
###### Observation 2.5.
If $`B`$ is a ucus Banach space and $`\varrho :G𝐎(B)`$, then the duality map $`:\mathrm{S}(B)\mathrm{S}(B^{})`$ between the unit spheres intertwines the actions of $`\varrho (G)`$ and $`\varrho ^{}(G)`$. In particular it maps the set of $`\varrho (G)`$-fixed unit vectors to the set of $`\varrho ^{}(G)`$-fixed unit vectors.
###### Proposition 2.6.
Let $`\varrho `$ be a uniformly equicontinuous linear representation of $`G`$ on a superreflexive space $`B`$, let $`B^{\varrho (G)}`$ denote the subspace of $`\varrho (G)`$-fixed vectors in $`B`$, and let $`B^{}=B^{}(\varrho )`$ be the annihilator of $`(B^{})^{\varrho ^{}(G)}`$ in $`B`$. Then
$$B=B^{\varrho (G)}B^{}(\varrho ).$$
Furthermore, the decomposition is canonical in the following sense: If we denote by $`p(\varrho )`$ and $`p^{}(\varrho )`$ the associated projections, then for every morphism of uniformly equicontinuous linear representations $`\varphi :(B_1,\varrho _1)(B_2,\varrho _2)`$, the following diagrams are commutative:
(2.iii)
$$\begin{array}{ccc}B_1& \stackrel{\varphi }{}& B_2\\ p\left(\varrho _1\right)& & p\left(\varrho _2\right)& & \\ B_1& \stackrel{\varphi }{}& B_2\end{array}\begin{array}{ccc}B_1& \stackrel{\varphi }{}& B_2\\ p^{}\left(\varrho _1\right)& & p^{}\left(\varrho _2\right)& & \\ B_1& \stackrel{\varphi }{}& B_2\end{array}$$
###### Remark 2.7.
The conclusion fails if we drop the superreflexivity assumption, see Example 2.29.
###### Proof of the proposition.
Choose a $`G`$-invariant uniformly convex and uniformly smooth norm on $`B`$, and the dual one on $`B^{}`$ (Proposition 2.3). For any unit vector $`xB^{\varrho (G)}`$ and arbitrary $`yB^{}`$
$$1=x,x^{}=xy,x^{}xyx^{}=xy.$$
Thus $`B^{\varrho (G)}B^{}=\{0\}`$ and $`B^{\varrho (G)}B^{}`$ is a closed subspace in $`B`$. It is also dense in $`B`$. Indeed if $`\lambda B^{}`$ is a unit vector vanishing on $`B^{}`$ it cannot vanish on $`B^{\varrho (G)}`$, because $`\lambda (B^{})^{\varrho ^{}(G)}`$ by the Hahn–Banach theorem, and hence $`\lambda ^{}B^{\varrho (G)}`$ and $`\lambda ^{},\lambda =1`$. Thus $`B^{\varrho (G)}B^{}=B`$.
The last assertion follows from the fact that $`\varphi (B_1^{\varrho _1})B_2^{\varrho _2}`$, and $`\varphi ^{}((B_2^{})^{\varrho _2})(B_1^{})^{\varrho _1}`$ yields $`\varphi (B_2^{})B_1^{}`$. ∎
###### Corollary 2.8.
The decomposition $`B=B^{\varrho (G)}B^{}`$ is preserved by the normalizer of $`\varrho (G)`$ in $`\mathrm{𝐆𝐋}(B)`$.
###### Corollary 2.9.
Let $`G=G_1\times G_2`$ be any product of two groups and $`B`$ a superreflexive space with a uniformly equicontinuous linear $`G`$-representation $`\varrho `$. Then there is a canonical $`G`$-invariant decomposition
$$B=B^{\varrho (G)}B_0B_1B_2$$
such that $`B^{\varrho (G_i)}=B^{\varrho (G)}B_i`$ for $`i=1,2`$.
###### Proposition 2.10.
Let $`\varrho `$ be a uniformly equicontinuous linear $`G`$-representation on a superreflexive space $`B`$. Then
1. $`B^{\varrho (G)}`$ is isomorphic to $`B/B^{}`$ as topological vector spaces.
2. $`B^{}`$ is isomorphic to $`B/B^{\varrho (G)}`$ as $`G`$-representations.
3. $`(B^{\varrho (G)})^{}`$ is isomorphic to $`B^{}/(B^{})^{}`$ as topological vector spaces.
4. $`(B^{})^{}`$ is isomorphic to $`(B^{})^{}`$ as $`G`$-representations.
5. $`B^{}`$ almost has invariants if and only if $`(B^{})^{}`$ almost has invariants.
6. If $`0ABC0`$ is an exact sequence of uniformly equicontinuous linear $`G`$-representations on superreflexive spaces, then $`B^{}`$ almost has invariant vectors if and only if $`A^{}`$ or $`C^{}`$ do.
If $`B`$ is equipped with a compatible uniformly convex and uniformly smooth $`G`$-invariant norm, then the natural isomorphisms in (1) and (3) are isometric.
###### Proof.
Equip $`B`$ with a $`G`$-invariant ucus norm (Proposition 2.3).
By the open mapping theorem the maps $`p:BB^{\varrho (G)}`$ and $`p^{}:BB^{}`$ induce isomorphisms of topological vector spaces
$$(1)\stackrel{~}{p}:B/B^{}B^{\varrho (G)},(2)\stackrel{~}{p^{}}:B/B^{\varrho (G)}B^{}.$$
By 2.5 $`(B^{\varrho (G)})^{}`$ is $`(B^{})^{\varrho ^{}(G)}`$ and the latter is isomorphic to $`B^{}/(B^{})^{}`$. This proves (3).
To see that (1) and (3) are isometric (with respect to the norms corresponding to any ucus $`G`$-invariant norm on $`B`$) we note that the isomorphisms above satisfy $`\stackrel{~}{p}^1,\stackrel{~}{p^{}}^11`$ by the definition of the norm on a quotient space. Furthermore, for $`v\mathrm{S}(B^{\varrho (G)})`$, we have $`v^{}\mathrm{S}((B^{})^{\varrho ^{}(G)})`$, hence
$`\stackrel{~}{p}^1(v)_{B/B^{}}`$ $`=`$ $`inf\{v+v^{}_B:v^{}B^{}\}`$
$``$ $`inf\{v+v^{},v^{}:v^{}B^{}\}=v,v^{}=1.`$
Hence $`\stackrel{~}{p}`$ is an isometry $`B/B^{}B^{\varrho (G)}`$. Similarly, $`(B^{\varrho (G)})^{}=(B^{})^{\varrho ^{}(G)}B^{}/(B^{})^{}`$.
In general Banach spaces the dual $`E^{}`$ of a subspace $`E<F`$ is isometric to the quotient $`F^{}/E^{}`$ by the annihilator $`E^{}<F^{}`$ of $`E`$. Thus with respect to a ucus norm on $`B`$ and the above spaces, $`(B^{})^{}`$ is isometric to $`B^{}/(B^{})^{\varrho ^{}(G)}`$ as Banach spaces, while the latter is isomorphic to $`(B^{})^{}`$ as a topological vector space by (2). Whence (4).
(5) Assume that there exist $`x_n\mathrm{S}(B^{})`$ with $`\mathrm{diam}(\varrho (K)x_n)0`$. The uniformly continuous map $`:\mathrm{S}(B)\mathrm{S}(B^{})`$ takes vectors $`x_n\mathrm{S}(B^{})`$ to vectors $`x_n^{}\mathrm{S}(B^{})`$ with $`\mathrm{diam}(\varrho ^{}(K)x_n^{})0`$. Since $`\{x_n^{}\}`$ are uniformly separated from $`(B^{})^{\varrho ^{}(G)}`$ their normalized projection $`y_n^{}`$ to $`(B^{})^{}`$ still satisfy $`\mathrm{diam}(\varrho ^{}(K)y_n^{})0`$.
(6) As $`A^{}`$ maps into $`B^{}`$, if $`A^{}`$ almost has invariants, then so does $`B^{}`$. If $`C^{}`$ almost has invariants, then so does $`(C^{})^{}`$, hence $`(B^{})^{}`$, hence $`B^{}`$. On the other hand, assume $`B^{}`$ almost has invariant unit vectors $`b_n`$. Assume for simplicity that $`A=A^{}`$, $`B=B^{}`$ and $`C=C^{}`$. Note that $`C`$ is isomorphic to $`B/A`$, and denote by $`\pi :BC`$ the projection. Then either $`\pi (b_n)`$ converges to $`0C`$, then there exist $`a_n`$ such that $`b_na_n`$ converges to $`0B`$, and the normalized sequence $`(\frac{a_n}{a_n})`$ is almost invariant in $`A`$, or there exist a subsequence $`b_{n_k}`$ with $`inf_k\pi (b_{n_k})>0`$, and then the normalized sequence $`(\frac{\pi (b_{n_k})}{\pi (b_{n_k})})`$ is almost invariant in $`C`$. ∎
###### Remark 2.11.
For ucus Banach space $`B`$, Definition 1.1 of property $`(T_B)`$ can be rephrased as follows: For any representation $`\varrho :G𝐎(B)`$, the restriction $`\varrho ^{}:G𝐎(B^{})`$ of $`\varrho `$ to the invariant subspace $`B^{}`$ complement to $`B^{\varrho (G)}`$ does not almost have invariant vectors, i.e. for some compact $`KG`$ and $`\epsilon >0`$
$$v\mathrm{S}(B^{})gK\text{s.t.}\varrho (g)vv\epsilon .$$
Hence item (4) gives:
###### Corollary 2.12.
Let $`B`$ be a ucus Banach space, and $`G`$ be a locally compact group. Then $`G`$ has property $`(T_B)`$ iff it has $`(T_B^{})`$.
### 2.d. Affine Actions
The affine group $`\mathrm{Aff}(V)`$ of a real affine space $`V`$ (a vector space who forgot its origin) consists of invertible maps satisfying:
$$T(tx+(1t)y)=tT(x)+(1t)T(y),(t𝐑,x,yV)$$
The group $`\mathrm{Aff}(V)`$ is a semi-direct product $`\mathrm{Aff}(V)=\mathrm{𝐆𝐋}(V)V`$, i.e. an invertible affine map $`T`$ has the form $`T(x)=Lx+b`$ where $`L\mathrm{𝐆𝐋}(V)`$ is linear invertible.
An affine action of a group $`G`$ on $`V`$, i.e. a homomorphism $`G\mathrm{Aff}(V)`$, has the form
$$gx=\varrho (g)x+c(g),$$
where $`\varrho :G\mathrm{𝐆𝐋}(V)`$ is a linear $`G`$-representation (we call it the *linear part* of the action) and $`c:GB`$ is a $`\varrho `$-*cocycle*, namely an element of the Abelian group
(2.iv)
$$Z^1(\varrho )=\{c:GV:c(gh)=\varrho (g)c(h)+c(g),g,hG\}.$$
The group $`Z^1(\varrho )`$ of $`\varrho `$-cocycles contains the subgroup of $`\varrho `$-*coboundaries*
(2.v)
$$B^1(\varrho )=\{c(g)=v\varrho (g)v:vV\}.$$
$`Z^1(\varrho )`$ describes all affine $`G`$-actions on $`V`$ with linear part $`\varrho `$, and $`B^1(\varrho )`$ corresponds to those affine actions which have a $`G`$-fixed point (namely $`v`$ in (2.v)). This description involves the choice of reference point – the origin – in the space. Two cocycles differing by a coboundary can be though of defining the same affine action viewed from different reference points. The first cohomology of $`G`$ with $`\varrho `$-coefficients is the Abelian group
$$H^1(G,\varrho )=Z^1(\varrho )/B^1(\varrho ).$$
It describes different types of actions in the above sense. $`H^1(G,\varrho )=0`$ iff any affine $`G`$-action on $`V`$ with linear part $`\varrho `$ has a fixed point.
For a Banach space $`B`$ denote by $`\mathrm{Isom}(B)`$ the group of isometries of $`B`$ as a metric space. It is a classical theorem of Mazur–Ulam that any surjective isometry $`T`$ of a (real) Banach space $`B`$ is necessarily affine $`T(x)=Lx+c`$ with linear part $`L𝐎(B)`$ being isometric. (This theorem is elementary when $`B`$ is strictly convex; comopare Lemma 6.1). Hence $`\mathrm{Isom}(B)=𝐎(B)B`$.
Now suppose that $`V`$ is a superreflexive topological vector space. Recall that a group $`G`$ of affine self maps is uniformly equicontinuous if it satisfies (2.ii). This condition is equivalent to uniform equicontinuity of the linear part $`\varrho :G\mathrm{𝐆𝐋}(V)`$.
###### Proposition 2.13.
For a superreflexive topological vector space $`V`$ and a group of transformations $`G`$ of $`V`$ the following conditions are equivalent:
1. $`G`$ is uniformly equicontinuous group of affine transformations of $`V`$.
2. $`G`$ acts by uniformly Lipschitz affine transformations with respect to any/all compatible norms on $`V`$.
3. $`G`$ acts by affine isometries with respect to some compatible norm on $`V`$.
4. $`G`$ acts by affine isometries with respect to some uniformly convex and uniformly smooth compatible norm on $`V`$.
###### Proof.
Apply Proposition 2.3 to the linear part of the affine action, using Mazur–Ulam to deduce in (3) that the action is affine. ∎
If $`G`$ is a topological group acting by affine transformations on a topological vector space $`V`$, continuity of the action
$$G\times VV,gx=\varrho (g)x+c(g)$$
is equivalent to continuity of the linear part $`G\times VV`$ and the continuity of the cocycle $`c:GV`$. Indeed $`c(g)=g0`$, and $`\varrho (g)x=gxc(g)`$.
Hence in the context of topological groups, affine actions should be assumed continuous, and $`Z^1(G,\varrho )`$ will include only continuous cocycles $`c:GV`$ (we assume that the linear part $`\varrho `$ is continuous as well). If $`G`$ is a locally compact $`\sigma `$-compact group, then $`Z^1(\varrho )`$ has a natural structure of a Fréchet space with respect to the family of semi-norms
$$c_K=\underset{gK}{sup}c(g)_V$$
where $`KG`$ runs over a countable family of compact subsets which cover $`G`$ and $`_V`$ is a norm inducing the topology of $`V`$. Moreover, if $`G`$ is compactly generated (e.g. if $`G`$ has property $`(T)`$) say by $`K_0`$, then $`c_{K_0}`$ is a norm on $`Z^1(\varrho )`$ (note that any cocycle $`cZ^1(\varrho )`$ is completely determined by its values on a generating set), and $`Z^1(\varrho )`$ is a Banach space with respect to this norm. We remark that in general $`B^1(\varrho )`$ is not closed in $`Z^1(\varrho )`$ (this is the idea behind the $`(F_B)(T_B)`$ argument of Guichardet – see Section 3).
###### Lemma 2.14.
For a uniform equicontinuous affine action of a group $`G`$ on a superreflexive space $`B`$, the following are equivalent:
1. There exists a bounded $`G`$-orbit.
2. All $`G`$-orbits are bounded
3. $`G`$ fixes a point in $`B`$.
4. $`G`$ preserves a (Borel regular) probability measure on $`B`$.
Note that the notion of a subset $`EV`$ being bounded, means that for any open neighbourhood $`U`$ of $`0V`$ there is some $`t𝐑`$ so that $`EtU`$. This notion agrees with the notion of being bounded with respect to any compatible norm on $`V`$.
###### Proof.
Introduce a $`G`$-invariant uniformly convex norm on $`V`$ (Proposition 2.13). The only non-trivial implications are $`(4)(1)(3)`$. For the first, let $`\mu `$ be a $`G`$-invariant probability on $`B`$. Since $`B`$ is a countable union of closed bounded sets, there is a closed bounded set $`AB`$ with $`\mu (A)>1/2`$. For all $`gG`$ we have $`\mu (gA)>1/2`$ hence $`gAA\mathrm{}`$. It follows that the $`G`$-orbit of every point of $`A`$ is bounded.
The latter implication follows by considering the circumcentre (compare Section 2.a) of the given bounded $`G`$-orbit. ∎
###### Proposition 2.15.
Let $`B`$ be a ucus Banach space. Then
1. Any finite (or compact) group has properties $`(T_B)`$ and $`(F_B)`$.
2. Properties $`(T_B)`$ and $`(F_B)`$ pass to quotient groups.
3. If $`G=G_1\times \mathrm{}\times G_n`$ is a finite product of topological groups then $`G`$ has property $`(T_B)`$ (resp. $`(F_B)`$) iff all $`G_i`$ have this property.
###### Proof.
(1) and (2) are straightforward, (3) follows from Corollary 2.9. ∎
### 2.e. Special Properties of $`L^p(\mu )`$-Spaces
In this section we collect some special properties of the Banach spaces $`L^p(\mu )`$ which will be used in the proofs.
By an $`L^p(\mu )`$, or $`L^p(X,\mu )`$ space we mean the usual space of equivalence classes (modulo null sets) of measurable $`p`$-integrable functions $`f:X𝐑`$, where $`\mu `$ is a positive $`\sigma `$-finite measure defined on a standard Borel space $`(X,)`$. If $`1<p<\mathrm{}`$ then $`L^p(\mu )`$ is ucus, whilst $`L^1(\mu )`$ and $`L^{\mathrm{}}(\mu )`$ are not (they are not even strictly convex). For $`1p<\mathrm{}`$ the dual to $`L^p(\mu )`$ is $`L^q(\mu )`$ where $`1<q\mathrm{}`$ is determined by $`q=p/(p1)`$.
The space $`L^p([0,1],\text{Lebesgue})`$ is usually denoted by $`L^p`$. Any $`L^p(\mu )`$-space with *non-atomic* finite or $`\sigma `$-finite measure $`\mu `$ is isometrically isomorphic to $`L^p`$. Indeed let $`\phi L^1(\mu )`$ be a strictly positive measurable function with integral one and let $`\mu _1`$ be given by $`d\mu _1=\phi d\mu `$. Then
$$fL^p(\mu )f\phi ^{1/p}L^p(\mu _1)$$
is a surjective isometry. Since any non-atomic standard probability spaces is isomorphic to $`[0,1]`$ as a measure space, $`L^p(\mu _1)L^p`$. If $`\mu `$ is purely atomic then a similar argument gives an isomorphism of $`L^p(\mu )`$ with finite or infinite dimensional $`\mathrm{}^p`$ space. A general $`L^p(\mu )`$ space is therefore isometrically isomorphic to a direct sum of $`L^p`$ and $`\mathrm{}^p`$ components.
More generally, for another Banach space $`B`$, one defines the spaces $`L^p(\mu ,B)`$ of $`B`$-valued function classes by means of the Bochner integral. We refer the reader to \[DU\] for details; we recall here that the dual of $`L^p(\mu ,B)`$ is $`L^q(\mu ,B^{})`$ through the natural pairing for all $`1p<\mathrm{}`$, but only when $`B`$ has the Radon–Nikodým property – this includes all ucus spaces (see again \[DU\]). These spaces will be used in Section 8.b in order to *induce* isometric (linear or affine) actions.
Banach \[Ba\] and Lamperti \[Li\] (see also \[FJ, Theorem 3.25\]) classified the linear isometries of $`L^p(\mu )`$ as follows.
###### Theorem 2.16 (Banach, Lamperti).
For any $`1<p<\mathrm{}`$ where $`p2`$, any linear isometry $`U`$ of $`L^p(X,,\mu )`$ has the form
$$Uf(x)=f(T(x))h(x)\left(\frac{dT_{}\mu }{d\mu }(x)\right)^{\frac{1}{p}}$$
where $`T`$ is a measurable, measure class preserving map of $`(X,\mu )`$, and $`h`$ is a measurable function with $`|h(x)|=1`$ a.e.
Let $`\mu =\mu _a+\mu _c`$ be the decomposition of $`\mu `$ into its atomic and continuous parts ($`\mu _a=\mu |_A`$ where $`AX`$ is the (at most countable) set of atoms of $`\mu `$). Then
$$L^p(\mu )=L^p(\mu _c)L^p(\mu _a)L^p\mathrm{}^p(A)\text{or just}\mathrm{}^p(A),$$
the latter case occurs if $`\mu =\mu _a`$ is a purely atomic measure. Note that it follows from Banach–Lamperti theorem that this decomposition is preserved by any linear isometry of $`L^p(\mu )`$. As $`\mathrm{}^p(A)`$ has a much smaller group of linear (or affine) isometries than $`L^p`$ we could restrict our attention only to the latter. However we shall not make use of this “simplification“.
Another useful tool in the study of $`L^p`$-spaces is the Mazur map.
###### Theorem 2.17 ( \[BL, Theorem 9.1\] ).
Let $`\mu `$ be a $`\sigma `$-finite measure. For any $`1p,q<\mathrm{}`$ the *Mazur map* $`M_{p,q}:L^p(\mu )L^q(\mu )`$ defined by
$$M_{p,q}(f)=\mathrm{sign}(f)|f|^{\frac{p}{q}}$$
is a (non-linear) map which induces a uniformly continuous homeomorphism between the unit spheres $`M_{p,q}:\mathrm{S}(L^p(\mu ))\mathrm{S}(L^q(\mu ))`$.
(Note that if $`p,q1`$ and $`p^1+q^1=1`$ then the restriction of $`M_{p,q}`$ to the unit spheres is just the duality map $`:\mathrm{S}(L^p(\mu ))\mathrm{S}(L^p(\mu )^{})`$).
In the proofs of Theorems A and B, the results for subspaces and quotients are deduced from the $`L^p(\mu )`$ case using the following theorem of Hardin about extension of isometries defined on subspaces of $`L^p(\mu )`$. The formulation we give here is not quite identical to the original, but it easily follows from it and from its proof (see \[Ha, Theorem 4.2\] or \[FJ, Theorem 3.3.14\]).
###### Theorem 2.18 (Hardin).
Let $`(X,,\mu )`$ be a measure space. For every closed subspace $`FL^p(X,\mu )`$, there is a canonical extension $`F\stackrel{~}{F}L^p(\mu )`$ which is isometric to $`L^p(X^{},\mu ^{})`$ for some other measure space $`(X^{},\mu ^{})`$. Furthermore, if $`1<p2𝐙`$, then every linear isometry $`U:FL^p(Y,\nu )`$ extends uniquely to a surjective linear isometry $`\stackrel{~}{U}:\stackrel{~}{F}\stackrel{~}{UF}L^p(Y,\nu )`$.
###### Remark 2.19.
If $`^{}`$ is the minimal sub $`\sigma `$-algebra with respect to which all the functions in $`F`$ are measurable, then $`\stackrel{~}{F}=L^p(X,^{},\mu )`$.
A straightforward consequence is the following:
###### Corollary 2.20.
Let $`1<p2𝐙`$, and let $`FL^p(X,\mu )`$ be a closed subspace. Let $`\varrho `$ be a linear isometric representation of the group $`G`$ on $`F`$. Then there is some linear isometric $`G`$-representation $`\varrho ^{}`$ of $`G`$ on some other space $`L^p(X^{},\mu ^{})`$, and a linear $`G`$-equivariant isometric embedding $`FL^p(X^{},\mu ^{})`$.
Another important fact about $`L^p(\mu )`$-spaces, this time for $`p(0,2]`$, is that $`B=L^p(\mu )`$ has an embedding $`j:B`$ into the unit sphere of a Hilbert space so that $`j(x),j(y)=xy^p`$. Having such an embedding is equivalent (via the classical result of I.J. Schoenberg, see \[BHV\]) to the following:
###### Proposition 2.21.
For $`0<p<2`$ and any $`s>0`$ the function $`fe^{sf^p}`$ is positive definite on $`L^p(\mu )`$, i.e. for any finite collection $`f_iL^p(\mu )`$ and any $`\lambda _i𝐂`$
$$\underset{i,j}{}e^{sf_if_j^p}\lambda _i\overline{\lambda _j}0.$$
In fact, more is known: It was shown by Bretagnolle, Dacunha-Castelle and Krivine \[BDCK\] (c.f. \[WW, 5.1\]) that, for $`1p2`$, a Banach space $`X`$ is isometric to a closed subspace of $`L^p(\mu )`$ iff $`e^{s^p}`$ is a positive definite function on $`X`$ for any $`s>0`$.
### 2.f. Some Easy Counterexamples and Remarks
###### Example 2.22 ($`(T_B)\Rightarrow ̸(F_B)`$).
Let $`B`$ be a Banach space with $`𝐎(B)𝐙/2𝐙`$, i.e. a space where the only linear isometries are the identity and the antipodal map $`xx`$. A trivial example of such a space is the line $`B=𝐑`$, but it is not hard to construct such spaces of arbitrary dimensions even within the class of ucus Banach spaces (by considering e.g. sufficiently asymetric convex sets in Hilbert space and choosing the corresponding norm). Clearly for such a space any group has property $`(T_B)`$. However the groups $`𝐙`$ or $`𝐑`$ or any group $`G`$ with sufficiently large Abelianization $`G/[G,G]`$ would fail to have property $`(F_B)`$ for it would admit an isometric action by translations: $`nx:=x+nx_0`$ where $`0x_0B`$ is arbitrary. However groups with trivial Abelianization would also have property $`(F_B)`$ on such an asymmetric Banach space $`B`$.
###### Example 2.23 ($`(T)\Rightarrow ̸(F_B)`$).
Suppose $`G`$ is locally compact non-compact (e.g. $`G=\mathrm{𝐒𝐋}_3(𝐑)`$ or $`G=\mathrm{𝐒𝐋}_3(𝐙)`$). Fix a Haar measure on $`G`$ and let
$$B=L_0^1(G)=\{fL^1(G):f𝑑g=0\}$$
be the codimension one subspace of functions with $`0`$ mean. Then $`B`$ is isometric to the affine subspace $`\{fL^1(G):f=1\}`$ on which $`G`$ acts isometrically by translations without fixed points. Hence $`G`$ does not have property $`(F_B)`$. This Banach space is not ucus. Notice that in this example all orbits are bounded regardless of $`G`$.
###### Remark 2.24.
Haagerup and Przybyszewska \[HP\] showed that any locally compact group $`G`$ admits a proper isometric action on the strictly convex space $`_{n=1}^{\mathrm{}}L^{2n}(G)`$.
###### Example 2.25 ($`(T)\Rightarrow ̸(T_B)`$).
Let $`G`$ be as in Example 2.23. Consider the space $`B=C_0(G)`$ of continuous real valued functions on $`G`$ which tend to $`0`$ at $`\mathrm{}`$ with the sup ($`L^{\mathrm{}}`$-) norm. The action of $`G`$ on $`B`$ by translations is a linear isometric action. A function $`fB`$ which decays very slowly forms an “almost invariant vector”. On the other hand there are no non-zero invariant vectors. Hence $`G`$ does not have property $`(T_B)`$.
###### Remark 2.26.
Since any separable Banach space is a quotient of $`\mathrm{}^1`$, Example 2.25 shows that case (iii) of Theorem A cannot be extended to $`p=1`$.
###### Example 2.27 (Remarks 1.8(2)).
Let $`G=G_1\times G_2`$ be any product of non-compact locally compact groups (e.g. $`G=𝐙\times 𝐙`$). Let $`B=L_0^1(G)`$ as in Example 2.23. Then $`H^1(G,B)0`$, but there are no non-zero $`G_i`$-fixed vectors in the associated linear representation. Thus the product formula of Remarks 1.8(2) cannot hold for $`B`$.
Let us make some remarks about Kazhdan’s property $`(T)`$ and property $`(T_B)`$ as in 1.1 and 2.11. Given a unitary representation $`(\varrho ,)`$ of a locally compact group $`G`$, a compact subset $`KG`$ and $`\epsilon >0`$, one says that a vector $`0v`$ is $`(K,\epsilon )`$- *almost invariant* if
$$\underset{gK}{sup}\varrho (g)vv<\epsilon v.$$
A locally compact group $`G`$ has Kazhdan’s property $`(T)`$ if and only if it satisfies the following equivalent conditions:
1. For any unitary $`G`$-representation $`(\varrho ,)`$ there exists a compact $`KG`$ and an $`\epsilon >0`$ so that the $`G`$-representation $`\varrho ^{}`$ on $`(^{\varrho (G)})^{}/^{\varrho (G)}`$ has no $`(K,\epsilon )`$-almost invariant vectors.
2. There exist a compact $`KG`$ and an $`\epsilon >0`$ so that all non-trivial irreducible unitary $`G`$-representations $`(\varrho ,)`$ have no $`(K,\epsilon )`$-almost invariant vectors.
3. There exist a compact $`KG`$ and an $`\epsilon >0`$ so that for all unitary $`G`$-representations $`(\varrho ,)`$ the $`G`$-representation $`\varrho ^{}`$ on $`(^{\varrho (G)})^{}/^{\varrho (G)}`$ has no $`(K,\epsilon )`$-almost invariant vectors.
In the above, (3) clearly implies both (1) and (2). In showing (1)$``$(3) one uses the fact that the category of Hilbert spaces and unitary representations is closed under $`\mathrm{}^2`$ sums and $`L^2`$-integration. The fact that any unitary representation decomposes as an $`L^2`$-integral of irreducible ones gives (2)$``$(3).
###### Remark 2.28.
Definition 1.1 (Remark 2.11) of property $`(T_B)`$ is modeled on (1) above. There does not seem to be any reasonable theory of irreducible representations (and decomposition into irreducibles) for Banach spaces other than Hilbert ones. Hence form (2) of property $`(T)`$ does not seem to have a Banach space generalization. As for (3), for any given $`1<p<\mathrm{}`$ the class of $`L^p(\mu )`$-spaces is closed under taking $`\mathrm{}^p`$-sums (and $`L^p`$-integrals) and hence for groups with property $`(T_{L^p})`$ an analogue of (3) holds, namely there exist $`KG`$ and $`\epsilon >0`$ which are good for all $`\varrho :G𝐎(L^p)`$. Also, if a group $`G`$ has property $`(T_B)`$ for *all* ucus Banach spaces $`B`$ (conjecturally all higher rank groups and their lattices) then for every ucus Banach space $`B`$ there is $`(K,\epsilon )`$ which is good for all linear isometric representations $`G𝐎(B)`$. This uses the fact that $`L^2(\mu ,B)`$ is a ucus if $`B`$ is (see Lemma 8.6 below).
Finally, we justify Remark 2.7:
###### Example 2.29.
Let $`G`$ be a discrete group and consider the Banach space $`B=\mathrm{}^{\mathrm{}}(G)`$ with the (linear isometric) regular $`G`$-representation $`\varrho `$. Then one shows that the space $`B^{\varrho (G)}`$ (which consists of the constant functions) admits a $`G`$-invariant complement (if and) *only if $`G`$ is amenable*. Indeed, the Riesz space (or Banach lattice) structure of $`B`$ allows to take the “absolute value” of any linear functional on $`B`$; renormalizing the absolute value of any non-zero invariant functional would yield an invariant mean on $`G`$. Alternatively, one can argue similarly on the Banach space of continuous functions on any compact topological $`G`$-space.
We point out that nevertheless the space $`B^{}`$ is well-defined for any topological vector $`G`$-space $`B`$; in the case at hand, we have $`B^{}=B`$ which shows why it cannot be a complement for $`B^{\varrho (G)}0`$.
## 3. Proof of Theorem 1.3
### 3.a. Guichardet: $`(F_B)(T_B)`$
###### Proof.
Assume $`G`$ does not have $`(T_E)`$, where $`E`$ is a Banach space, and let $`\varrho :G𝐎(E)`$ be a representation such that $`E/E^{\varrho (G)}`$ admits almost invariant vectors. In order to show that $`H^1(G,\varrho )\{0\}`$ it suffices to prove that $`B^1(G,\varrho )Z^1(G,\varrho )`$ is not closed.
As was mentioned in Section 2 the space of $`\varrho `$-cocycles $`Z^1(G,\varrho )`$ is always a Fréchet space (and even a Banach space if $`G`$ is compactly generated). Note that $`B^1(G,\varrho )`$ is the image of the bounded linear map
$$\tau :EZ^1(G,\varrho ),\left(\tau (v)\right)(g)=v\varrho ^{}(g)v.$$
If $`\tau (E)`$ were closed, and hence a Fréchet space, the open mapping theorem would imply that $`\tau ^1:B^1(G,\varrho )E/E^{\varrho (G)}`$ is a bounded map. That would mean that for some $`M<\mathrm{}`$ and a compact $`KG`$
$$vM\tau (v)_K=M\underset{gK}{sup}\varrho (g)vv,vE/E^{\varrho (G)}$$
contrary to the assumption that $`\varrho `$ almost contains invariant vectors. ∎
### 3.b. $`(T)(F_{L^p})`$, $`0<p2`$
###### Proof.
Let $`G`$ be a locally compact group with Kazhdan’s property $`(T)`$ acting by affine isometries on a closed subspace $`BL^p(\mu )`$ with $`0<p2`$. Using Proposition 2.21 and a slight modification of a Delorme–Guichardet argument for $`(T)(FH)`$ we shall prove that such an action has bounded orbits. For $`1<p2`$ uniform convexity of $`BL^p(\mu )`$ yields a $`G`$ fixed point using Lemma 2.14.
Proposition 2.21 allows to define a family, indexed by $`s>0`$, of Hilbert space $`_s`$, embeddings $`U_s:B\mathrm{S}(_s)`$ and unitary representations $`\pi _s:G𝐎(_s)`$ with the following properties:
1. The image $`U_s(B)`$ spans a dense subspace of $`_s`$;
2. $`U_s(x),U_s(y)=e^{sxy^p}`$ for all $`x,yB`$;
3. $`U_s(gx)=\pi _s(g)U_s(x)`$ for all $`xB`$, $`gG`$.
Indeed, one constructs $`_s`$ as the completion of the pre-Hilbert space whose vectors are finite linear combinations $`a_ix_i`$ of points $`x_iB`$, and the inner product is given by
$$a_ix_i,b_jy_j=\underset{i,j}{}a_i\overline{b_j}e^{sx_iy_j^p}.$$
The representation $`\pi _s`$ can be constructed (and is uniquely determined) by property (3).
Since $`G`$ is assumed to have Kazhdan’s property $`(T)`$, for some compact subset $`KG`$ and $`\epsilon >0`$, any unitary $`G`$-representation with $`(K,\epsilon )`$-almost invariant vectors has a non-trivial invariant vector.
Let $`x_0B`$ be fixed. The isometric $`G`$-action is continuous, so $`Kx_0`$ is a compact and hence bounded subset of $`B`$, hence:
$$R_0=\underset{gK}{sup}gx_0x_0<\mathrm{}.$$
For the unit vectors $`u_s=U_s(x_0)\mathrm{S}(_s)`$ we have
$$\underset{gK}{\mathrm{min}}|\pi _s(g)u_s,u_s|e^{sR_0^p}1\text{as}s0.$$
In particular for a sufficiently small $`s>0`$, $`\mathrm{max}_{gK}\pi _s(g)u_su_s<\epsilon `$. Let us fix such an $`s`$, and rely on property $`(T)`$ to deduce that $`\pi _s`$ has an invariant vector $`v\mathrm{S}(_s)`$.
We claim that $`G`$ must have bounded orbits for its affine isometric action on $`B`$. Indeed, otherwise there would exist a sequence $`g_nG`$ so that
$$g_nxy\mathrm{}\text{and hence}\pi _s(g_n)U_s(x),U_s(y)0$$
for all $`x,yB`$. This implies that $`\pi _s(g_n)w,u0`$ for any $`w,u\text{span}(U_s(B))`$, and since $`\text{span}(U_s(B))`$ is dense in $`_s`$, for any $`w,u_s`$. Taking $`w=u=v`$, we get a contradiction. Therefore the affine isometric $`G`$-action on $`B`$ has bounded orbits, and hence fixes a point in case of $`1<p2`$. ∎
### 3.c. Fisher–Margulis: $`(T)(F_{L^p})`$, $`p<2+\epsilon (G)`$
Let $`G`$ have Kazhdan’s property $`(T)`$. Fix a compact generating subset $`K`$ of $`G`$.
###### Lemma 3.1.
There exists a constant $`C<\mathrm{}`$ and $`\epsilon >0`$ such that for any $`G`$-action by affine isometries on a closed subspace $`BL^p(\mu )`$ with $`p(2\epsilon ,2+\epsilon )`$ and any $`xB`$ there exists a point $`yB`$ with
$$xyC\mathrm{diam}(Kx),\mathrm{diam}(Ky)<\frac{\mathrm{diam}(Kx)}{2}.$$
###### Proof.
By contradiction there exists a sequence of subspaces $`B_nL^{p_n}`$ with $`p_n2`$, affine isometric $`G`$-actions on $`B_n`$ and points $`x_nB_n`$ so that, after a rescaling to achieve $`\mathrm{diam}(Kx_n)=1`$, we have
(3.i)
$$\mathrm{diam}(Ky)\frac{1}{2}y\mathrm{B}(x_n,n).$$
Passing to an ultraproduct of the spaces $`B_n`$ with the marked points $`x_n`$ and the corresponding $`G`$-actions, one obtains an isometric (hence also affine) $`G`$-action on a Hilbert space $``$, because the limit of $`L^p`$-parallelogram as $`p2`$ is the parallelogram identity, which characterizes Hilbert spaces. (The action is well-defined because $`K`$ generates $`G`$ and we ensured $`\mathrm{diam}(Kx_n)=1`$.) If $`G`$ is a topological group, one needs to ensure continuity of the limit action by selecting uniformly $`K`$-equicontinuous sets of vectors (as in \[Sh, 6.3\]; compare also \[CCS\]). Due to (3.i) this $`G`$-action has no fixed points, contradicting property $`(FH)`$ and hence $`(T)`$ of $`G`$. ∎
###### Proof of $`(F_B)`$ for $`BL^p(\mu )`$, $`2p<2+\epsilon (G)`$.
Now consider an arbitrary affine isometric $`G`$-action on a closed subspace $`BL^p`$ with $`|p2|<\epsilon `$ where $`\epsilon =\epsilon (G)>0`$ is as in the lemma. Define a sequence $`x_nB`$ inductively, starting from an arbitrary $`x_0`$. Given $`x_n`$, let $`R_n=\mathrm{diam}(Kx_n)`$. Then applying the lemma there exists $`x_{n+1}`$ within the ball $`\mathrm{B}(x_n,CR_n)`$ so that
$$\mathrm{diam}(Kx_{n+1})<R_n/2.$$
We get $`R_n<R_0/2^n`$ and $`x_{n+1}x_n<\mathrm{}`$. The limit of the Cauchy sequence $`\{x_n\}`$ is a $`G`$-fixed point. ∎
###### Question 3.2.
For a given group $`G`$ with property $`(T)`$, what can be said about the following invariant?
$$p(G):=inf\{p:\text{ }G\text{ fails to have }(F_B)\text{ for some closed subspace }BL^p\text{ }\}.$$
For instance, Pansu’s aforementioned result \[Pa\] shows that $`p(G)4n+2`$ for $`G=\mathrm{𝐒𝐩}_{n,1}(𝐑)`$.
## 4. Proof of Theorem A
We start with the first assertion of the theorem: $`(T)(T_B)`$ for $`B`$ being an $`L^p`$-related space as in (i), (ii) or (iii) in the theorem. We first reduce to the case (i) where $`B=L^p(\mu )`$ with $`1p<\mathrm{}`$. Then using Corollary 2.20 of Hardin’s extension theorem, $`(T_{L^p(\mu )})`$ implies $`(T_B)`$ for subspaces $`BL^p(\mu )`$ where $`p4,6,\mathrm{}`$ as in (ii), and the duality argument (Corollary 2.12) gives the result for quotients of $`L^q(\mu )`$ with $`q4/3,6/5,\mathrm{}`$ as in (iii). Hence it suffices to prove $`(T)(T_{L^p(\mu )})`$ for $`1p<\mathrm{}`$. We give two proofs for this implication.
Let us note that our restriction on $`p`$ and $`q`$ when taking subspaces/quotients comes from our use of Hardin’s theorem.
###### Question 4.1.
Does property $`(T)`$ implies property $`(T_B)`$ for any closed subspace and any quotient $`B`$ of $`L^p(\mu )`$ for any $`1<p<\mathrm{}`$ ?
### 4.a. Property $`(T)`$ Implies $`(T_{L^p(\mu )})`$ – First Proof
###### Proof.
Assuming that a locally compact group $`G`$ fails to have property $`(T_{L^p(\mu )})`$ for some $`1p<\mathrm{}`$, we are going to show that $`G`$ does not have (T). We may and will assume $`p2`$; write $`B=L^p(\mu )`$ and $`=L^2(\mu )`$. Using Remark 2.11 there is a representation $`\varrho :G𝐎(B)`$ so that for the canonical complement $`B^{}`$ of $`B^{\varrho (G)}`$ the restriction $`\varrho ^{}:G𝐎(B^{})`$ almost has invariant vectors, i.e. there exist unit vectors $`v_n\mathrm{S}(B^{})`$ so that
$$f_n(g)=\varrho (g)v_nv_n$$
converges to $`0`$ uniformly on compact subsets of $`G`$.
We shall obtain a related unitary, or orthogonal, representation $`\pi :G𝐎()`$ using the following:
###### Lemma 4.2.
For $`p2`$, the conjugation $`UM_{p,2}UM_{2,p}`$ by the non-linear Mazur map sends $`𝐎(B)`$ to $`𝐎()`$.
###### Proof.
Follows from Banach–Lamperti description of $`𝐎(B)`$ (Theorem 2.16) by calculation. ∎
Let us then define $`\pi :G𝐎()`$ by $`\pi (g)=M_{p,2}\varrho (g)M_{2,p}`$. Note that $`M_{p,2}`$ maps $`B^{\varrho (G)}`$ onto $`^{\pi (G)}`$.
As $`\mathrm{S}(B^{})`$ is uniformly separated (in fact is at distance $`1`$) from $`B^{\varrho (G)}`$, the uniform continuity of the Mazur map (Theorem 2.17) implies that $`u_n=M_{p,2}(v_n)`$ is a sequence in $`\mathrm{S}()`$ such that $`\text{dist}(u_n,^{\pi (G)})\delta >0`$ and $`\phi _n(g)=\pi (g)u_nu_n0`$ uniformly on compact subsets of $`G`$. Let $`w_n`$ denote the projections of $`u_n`$ to $`=(^{\pi (G)})^{}`$. Then
$$w_n\delta >0\text{and}\pi (g)w_nw_n\phi _n(g)0$$
uniformly on compacta. Thus the restriction $`\pi ^{}`$ of $`\pi `$ to $`^{}`$ does not have $`G`$-invariant vectors, but almost does. Hence $`G`$ does not have Kazhdan’s property $`(T)`$. ∎
###### Remark 4.3.
In fact, the above proof has established the following more specific statement. Let $`G`$ act measurably on a $`\sigma `$-finite measure space. Denote by $`\varrho _p`$ the associated linear isometric representation on $`L^p`$, namely the quasi-regular representation twisted by the $`p`$-th root of the Radon–Nikodým derivative. Then, the existence of almost invariant vectors in $`L^p/(L^p)^{\varrho _p(G)}`$ is independent of $`1p<\mathrm{}`$.
### 4.b. Property $`(T)`$ Implies $`(T_{L^p(\mu )})`$ – Second Proof
###### Proof.
For $`1<p2`$ we have $`(T)(F_{L^p(\mu )})(T_{L^p(\mu )})`$ by Theorem 1.3 (1) and (2). Using duality (Corollary 2.12) this implication extends to $`L^p(\mu )`$ with $`2<p<\mathrm{}`$. ∎
### 4.c. Property $`(T_{L^p})`$ Implies $`(T)`$
###### Proof.
Assume that $`G`$ is not Kazhdan, i.e. $`G`$ admits a unitary representation $`\pi `$ almost containing (but not actually containing) non-trivial invariant vectors. Connes and Weiss \[CW\] showed how to find such a representation of the form $`L_0^2(\mu )`$. More precisely, they construct a measure-preserving, *ergodic* $`G`$-action on a probability space $`(X,\mu )`$ which admits a a sequence $`\{E_n\}`$ of *asymptotically invariant* measurable subsets, namely
(4.i)
$$gG\mu (gE_nE_n)0\text{whilst}\mu (E_n)=1/2.$$
Consider the unitary $`G`$-representation $`\pi ^{}`$ on $`L_0^2(\mu )`$ – the space of *zero mean* square integrable functions, which is the orthogonal complement of the constants. Then $`\pi ^{}`$ does not have non-trivial invariant vectors because of ergodicity; but it almost does, namely $`f_n=2\mathrm{𝟏}_{E_n}1`$.
For a given $`1p<\mathrm{}`$, consider the linear isometric $`G`$-representation $`\varrho `$ on $`B=L^p(\mu )`$, $`\varrho (g)f(x)=f(g^1x)`$. Then $`B^{\varrho (G)}=𝐑\mathbf{\hspace{0.17em}1}`$ – the constants, and its canonical complement is
$$B^{}=L_0^p(\mu )=\{fL^p(\mu ):f𝑑\mu =0\}.$$
The above sequence $`\{f_n\}`$ lies in $`L_0^p(\mu )`$, consists of unit vectors and still satisfies $`\varrho (g)f_nf_n_p0`$. Hence failing to have Kazhdan’s property $`(T)`$ a group $`G`$ does not have $`(T_{L^p(\mu )})`$ either.
In the original paper \[CW\], Connes and Weiss considered discrete groups. In a similar context the case of locally compact groups was also considered by Glasner and Weiss (see \[GW, Section 3\] and references therein). One way to treat the non-discrete case, is the following: start from a unitary representation $`\pi `$ of a given lcsc $`G`$ which has almost invariant vectors but no invariant ones, and apply the original Connes–Weiss Gaussian construction to the restriction $`\pi |_\mathrm{\Gamma }`$ of $`\pi `$ to some dense countable subgroup $`\mathrm{\Gamma }G`$. This gives an ergodic measure-preserving $`\mathrm{\Gamma }`$-action on a probability space $`(X,\mu )`$ with an asymptotically invariant sequence $`\{E_n\}`$ on $`X`$. The fact that the representation $`\pi |_\mathrm{\Gamma }`$ came from $`G`$ is manifested by the fact that it is continuous in the topology on $`\mathrm{\Gamma }`$ induced from $`G`$. It can be shown to imply that the $`\mathrm{\Gamma }`$-representation on $`L_0^2(X,\mu )`$ is also continuous, hence extends to $`G`$, and thus the $`\mathrm{\Gamma }`$-action on $`(X,\mu )`$ extends to a measurable $`G`$-action. This construction gives a uniform convergence in (4.i) on compact subsets of $`G`$. ∎
## 5. Fixed Point Property for Higher Rank Groups
### 5.a.
The objective of this section is to prove Theorem B; we start with some preliminaries for the *linear* part.
The first ingredient needed for the proof is an analogue of Howe–Moore’s theorem on vanishing of matrix coefficients, or rather its corollary analogous to Moore’s ergodicity theorem, extended to the framework of uniformly equicontinuous representations on superreflexive Banach spaces. The ucus Banach space version of Howe–Moore is due to Yehuda Shalom (unpublished). With his kind permission we have included the argument in Appendix 9. Here we shall use the following corollary, which we formulate for the case of simple groups.
###### Corollary 5.1 (Banach space analogue of Moore’s theorem).
Let $`k`$ be a local field and let $`G=𝐆(k)`$ be the $`k`$-points of a Zariski connected isotropic simple $`k`$-algebraic group $`𝐆`$. Let $`G^+`$ be the image of the simply connected form $`\stackrel{~}{G}`$ in $`G`$ under the cover map. Let $`HG^+`$ be a closed non-compact subgroup.
Then for any superreflexive space $`B`$ and any continuous uniformly equicontinuous linear $`G`$-representation $`\varrho :G^+\mathrm{𝐆𝐋}(B)`$, $`B^{\varrho (H)}=B^{\varrho (G^+)}`$ and the canonical complements with respect to both $`\varrho (G^+)`$ and $`\varrho (H)`$ coincide, and can be denoted just by $`B^{}`$.
###### Proof.
By Proposition 2.3, we may assume that $`B`$ is a ucus Banach space and $`\varrho `$ is a linear isometric representation $`\varrho :G𝐎(B)`$. Now the statement follows readily from Theorem 9.1. ∎
### 5.b.
The second ingredient is *strong relative property $`(T)`$*. It will be used to prove Claim 5.5 below which is the only part which is specific to $`L^p`$-like spaces. The rest of the argument applies to all affine isometric actions on ucus Banach spaces, or all uniformly equicontinuous affine actions on a superreflexive space.
###### Definition 5.2.
Let $`HU`$ be a semi-direct product of locally compact groups. We shall say that it has
if for any unitary representation $`\pi `$ of $`HU`$ for which $`H`$ almost has non-trivial invariant vectors, $`U`$ has invariant vectors.
where $`B`$ is a Banach space, if for any linear isometric representation $`\varrho :HU𝐎(B)`$ the linear isometric $`H`$-representation $`\varrho ^{}:H𝐎(B/B^{\varrho (U)})`$ does not almost have non-trivial invariant vectors.
###### Remarks 5.3.
1. The first definition is a variant of “relative property $`(T)`$”. The latter usually refers to a pair of groups $`G_0G`$ and requires that any unitary $`G`$-representation with $`G`$-almost invariant vectors, has non-trivial $`G_0`$-invariant vectors. Strong relative property $`(T)`$ for $`HU`$ implies, but is not equivalent to, relative property $`(T)`$ for $`(HU,U)`$. In fact $`\mathrm{𝐒𝐋}_2(𝐑)𝐑^2`$ has the strong relative $`(T)`$ and thus relative (T) as well, whilst its lattice $`\mathrm{𝐒𝐋}_2(𝐙)𝐙^2`$ does not have strong relative $`(T)`$ even though the pair $`(\mathrm{𝐒𝐋}_2(𝐙)𝐙^2,𝐙^2)`$ has relative property $`(T)`$. (For the latter, cf. M. Burger’s appendix in \[HV\]. For the former, consider the representation on $`\mathrm{}^2(𝐙^2)`$ induced by the affine action on $`𝐙^2`$.)
2. If $`B`$ is a ucus Banach space, then the canonical splitting with respect to $`\varrho (U)`$, namely $`B=B^{\varrho (U)}B^{}`$ is preserved by $`\varrho (H)`$ which normalizes $`\rho (U)`$ (Corollary 2.8). Hence, as in Remark 2.11, for ucus space $`B`$ strong relative property $`(T_B)`$ requires that the *restriction* of $`\varrho (H)`$ to $`B^{}`$ does not almost have invariant vectors. Strong relative $`(T_{})`$ for a Hilbert space $``$ is equivalent to the strong relative $`(T)`$.
###### Lemma 5.4.
A semi-direct product $`HU`$ with strong relative property $`(T)`$ has strong relative property $`(T_B)`$ for all $`L^p`$-related Banach spaces $`B`$ of types (i), (ii), (iii) as in Theorem A.
###### Proof.
This is analogous to the proof of $`(T)(T_B)`$ given in Section 4.a. First observe that the extension Theorem 2.20 and a duality argument (based on Proposition 2.10) reduce the statement to the case (i) of $`B=L^p(\mu )`$.
Thus we assume that $`B=L^p(\mu )`$ with $`p2`$, and $`\varrho :HU𝐎(B)`$ is a linear isometric representation. Let $`B=B^{\varrho (U)}B^{}`$ be the canonical splitting with respect to $`U`$. It is preserved by $`\varrho (H)`$ because $`H`$ normalizes $`U`$. Now let $`\pi =M_{p,2}\varrho M_{2,p}`$ be the conjugate of $`\varrho `$ by the Mazur map. Then $`\pi `$ is an orthogonal representation $`\pi :HU𝐎()`$ where $`=L^2(\mu )`$ (Lemma 4.2).
If $`HU`$ fails to have strong relative $`(T_B)`$, then there exist $`x_n\mathrm{S}(B^{})`$ so that $`\varrho (h)x_nx_n0`$ uniformly on compact subsets of $`H`$. Uniform continuity of $`M_{p,2}`$ and the fact that $`\text{dist}(\mathrm{S}(B^{}),\mathrm{S}(B^{\varrho (U)}))=1`$, imply that for $`v_n=M_{p,2}(x_n)`$
$$\text{dist}(v_n,^{\varrho (U)})\delta >0\pi (h)v_nv_n0$$
uniformly on compact subsets of $`H`$. Taking projections of $`v_n`$ to $`^{}`$ we show that in this case $`HU`$ does not have strong relative property $`(T)`$. ∎
### 5.c. Proof of Theorem B
We first show that we can assume that $`G`$ is connected and simply connected. Assuming that Theorem B is known for $`\stackrel{~}{G_0}`$ and lattices therein; we will prove it for $`G`$ and its lattices. For any affine isometric action of $`G`$ on $`B`$ there is an associated action of $`\stackrel{~}{G_0}`$, inflated via the covering map $`\stackrel{~}{G_0}G`$. $`\stackrel{~}{G_0}`$ has a fixed point by assumption, hence $`G`$ has a compact orbit, as the cokernel of the covering map is compact \[M5, Theorem 2.3.1(b)\]. It follows that $`G`$ has a fixed point as well. A similar argument applies to lattices: For a given lattice $`\mathrm{\Gamma }`$ in $`G`$ its $`\stackrel{~}{\mathrm{\Gamma }}`$ by the covering map is a lattice in $`\stackrel{~}{G_0}`$, and its projection is of finite index in $`\mathrm{\Gamma }`$. Every affine isometric action of $`\mathrm{\Gamma }`$gives rise to an affine isometric action of $`\stackrel{~}{\mathrm{\Gamma }}`$, which, by assumption, has a fixed point. It follows that $`\mathrm{\Gamma }`$ has a finite orbit, and therefore fixes a point.
Hereafter we will assume that $`G`$ is connected and simply connected. In that case $`G`$ decomposes into a direct product of simply connected almost simple groups $`G=G_i`$ \[M5, Proposition 1.4.10\].
In view of (the independent) Sections 8.a and 8.b, more specifically Proposition 8.8(2) and the discussion following Definition 8.2, property $`(F_B)`$ for $`G=G_i`$ is inherited by its lattices. Thus it suffices to consider the ambient group $`G=G_i`$ only. By Proposition 2.15(3) the statement reduces to that about almost-simple factors $`G_i`$.
So we are left proving the theorem for $`G=𝐆(k)`$, a higher rank connected, simply-connected, almost-simple group. Using Proposition 2.13, we assume that $`B`$ is a ucus Banach space and we consider a $`G`$-action on $`B`$ by affine isometries, with $`\varrho :G𝐎(B)`$ denoting the linear part of the action. Let $`B=B^{\varrho (G)}B^{}`$ be the canonical decomposition and $`\varrho ^{}:G𝐎(B^{})`$ denote the corresponding sub-representation.
###### Claim 5.5 (For $`L^p`$-like spaces).
$`G`$ contains a direct product $`A\times H`$ so that
1. The restriction $`\varrho ^{}|_H:H𝐎(B^{})`$ does not almost contain invariant vectors.
2. $`A`$ contains a non-trivial semisimple element, and in particular it is not compact.
###### Proof.
Any higher rank almost-simple group $`G=𝐆(k)`$ is known to contain a subgroup whose simply-connected cover is isomorphic to either $`G_0=\mathrm{𝐒𝐋}_3(k)`$ or $`G_0=\mathrm{𝐒𝐩}_4(k)`$ \[M5, Theorem 1.6.2\]. In the first case $`G_0=\mathrm{𝐒𝐋}_3(k)`$ contains the semi-direct product $`H_0U_0=\mathrm{𝐒𝐋}_2(k)k^2`$ embedded in $`\mathrm{𝐒𝐋}_3(k)`$ as
$$\left\{\left(\begin{array}{ccc}a& b& x\\ c& d& y\\ 0& 0& 1\end{array}\right)\right|adbc=1\}$$
where $`U_0k^2`$ is the subgroup given by $`a=d=1`$, $`b=c=0`$. It is normalized by the copy $`H_0`$ of $`\mathrm{𝐒𝐋}_2(k)`$ embedded in the upper left corner. Let $`A_0\mathrm{𝐒𝐋}_3(k)`$ be the subgroup $`\mathrm{diag}[\lambda ,\lambda ,\lambda ^2]`$, $`\lambda k^{}`$, which centralizes $`H_0`$ in $`G_0`$, and let $`A`$ and $`HU`$ denote the corresponding subgroups in $`G`$.
The semi-direct product $`\mathrm{𝐒𝐋}_2(k)k^2`$ is known to have strong relative property $`(T)`$. Hence it has strong relative property $`(T_B)`$ for $`L^p`$-related spaces $`B`$ (Lemma 5.4). By 5.1 we have $`B^{\varrho (G)}=B^{\varrho (U)}`$ and we have denoted by $`B^{}`$ the common canonical complement. Then (1) follows from the strong relative property $`(T_B)`$ for $`HU`$, while (2) is clear from the construction.
In the second case $`G`$ contains a copy of $`G_0=\mathrm{𝐒𝐩}_4(k)`$ which is usually defined as a subgroup of $`\mathrm{𝐒𝐋}_4(k)`$ by
$$\mathrm{𝐒𝐩}_4(k)=\{g\mathrm{𝐒𝐋}_4(k)^tgJg=J\},\text{where}J=\left(\begin{array}{cc}0& I\\ I& 0\end{array}\right).$$
The semi-direct product $`H_0U_0`$ embedded in $`\mathrm{𝐒𝐋}_4(k)`$ is
$$\left\{\left(\begin{array}{cc}A& B\\ 0& {}_{}{}^{t}A_{}^{1}\end{array}\right)\right|A\mathrm{𝐒𝐋}_2(k),^tB=A^1B(^tA)\}$$
with $`H_0`$ denoting the image $`A\mathrm{diag}[A,^tA^1]`$ of $`\mathrm{𝐒𝐋}_2(k)`$, and $`U_0`$ the normal Abelian subgroup
$$\{\left(\begin{array}{cc}I& B\\ 0& I\end{array}\right)|^tB=B\}.$$
The semi-direct product $`H_0U_0`$ actually lies in $`\mathrm{𝐒𝐩}_4(k)`$, it is isomorphic to $`\mathrm{𝐒𝐋}_2(k)S^2(k)`$, where $`S^2(k)`$ is the space of symmetric bilinear forms on $`k^2`$ with the natural $`\mathrm{𝐒𝐋}_2(k)`$ action. This semi-direct product is also known to have strong relative property $`(T)`$, and therefore strong relative $`(T_B)`$. $`H_0`$ is centralized by $`A_0=\{\mathrm{diag}[\lambda ,\lambda ,\lambda ^1,\lambda ^1]:\lambda k^{}\}`$. As in the $`G_0=\mathrm{𝐒𝐋}_3(k)`$ case, we conclude that the corresponding product $`A\times HG`$ satisfies (1) and (2). The claim is proved. ∎
We now turn to the affine isometric $`G`$-action defined by a $`\varrho `$-cocycle $`cZ^1(\varrho )`$. We shall prove that $`cB^1(\varrho )`$ i.e. that $`G`$ has a global fixed point. Write $`c(g)=c_0(g)+c^{}(g)`$ with $`c_0(g)B^{\varrho (G)}`$ and $`c^{}(g)B^{}`$ where $`B=B^{\varrho (G)}B^{}`$ is the canonical splitting. Then $`c_0:GB`$ is a homomorphism into the (additive) Abelian group. As $`G`$ has compact Abelianization, $`c_0(g)0`$, which means that the affine $`G`$-action preserves each affine subspace $`p+B^{}`$. Hence both the affine $`G`$-action and the representation can be restricted to $`B^{}`$.
Claim 5.5 provides an input for the following general lemma:
###### Lemma 5.6.
Let a direct product of topological groups $`A\times H`$ act by affine isometries on a Banach space $`B`$. Suppose that the associated linear isometric representation $`\varrho `$ restricted to $`H`$ does not almost have invariant vectors. Then the affine action of $`A`$ has bounded orbits in $`B`$. In particular, if $`B`$ is uniformly convex, then $`A`$ has a fixed point in $`B`$.
###### Remark 5.7.
In the uniformly convex case, this follows of course from the stronger splitting theorem (Theorem C); compare also with Theorem 7.1 below for the weaker assumption that the *product* does not almost have invariant vectors.
###### Proof of the lemma.
Let $`\varrho :A\times H𝐎(B)`$ and $`cZ^1(\varrho )`$ denote the associated linear isometric representation and the translation cocycle. The commutation relation between any $`hH`$ and $`aA`$ gives
$$c(h)+\varrho (h)c(a)=c(ha)=c(ah)=c(a)+\varrho (a)c(h)$$
which can be rewritten as
$$(I\varrho (h))c(a)=(I\varrho (a))c(h).$$
By the assumption on $`\varrho (H)`$, there exists a compact subset $`KH`$ and an $`\epsilon >0`$ so that $`\mathrm{max}_{hK}\varrho (h)vv\epsilon v`$ for all $`vB`$. Let $`R=\mathrm{max}_{hK}c(h)<\mathrm{}`$. Then for $`aA`$
$$\epsilon c(a)\underset{hK}{\mathrm{max}}(I\varrho (h))c(a)2R.$$
Hence $`sup_{aA}c(a)2R/\epsilon `$, i.e. the $`A`$-orbit of $`0`$ is bounded. If $`B`$ is uniformly convex then the circumcentre of this orbit is an $`A`$-fixed point as in Lemma 2.14. ∎
Let us restrict the $`G`$-action to $`B^{}`$. For $`gG`$ let $`\text{Fix}(g)`$ denote the set of $`g`$-fixed points in $`B^{}`$. It follows from Claim 5.5 and Lemma 5.6 that for some non-elliptic semisimple element $`aAG`$, $`\text{Fix}(a)B^{}`$ is non-empty.
At this point we propose two different ways to conclude the proof.
First argument. For any $`g`$ we define
$$U(g)=\{hG:\underset{n\mathrm{}}{lim}g^nhg^n=e\}.$$
In analogy to the *Mautner phenomenon* \[M5, II.3.2\], we remark that in any continuous $`G`$-action by isometries on any metric space, every $`g`$-fixed point is $`U(g)`$-fixed. Indeed, if $`gx=x`$ then $`d(hx,x)=d(g^nhg^nx,x)0`$. It follows that any $`g`$-fixed point is fixed by both $`U(g)`$ and $`U(g^1)`$. Back to the setting of Theorem B, we apply this to $`a`$ and find that any $`a`$-fixed point is fixed by $`G`$ because the latter is generated by $`U(a)U(a^1)`$, see \[M5, I.1.5.4(iii)\].
Second argument. Note that $`\text{Fix}(a)`$ cannot contain more than one point. Indeed, if $`x,y\text{Fix}(a)`$ then $`xyB^{}`$ is a fixed vector for the linear isometry $`\varrho (a)`$. Since the cyclic group $`a`$ is unbounded, the ucus analogue of Moore’s ergodicity 5.1 implies that $`B^{\varrho (g)}=B^{\varrho (G)}`$. Hence $`xy=0`$. Thus $`\text{Fix}(a)=\{x_0\}`$.
If $`g,hG`$ commute then $`\text{Fix}(g)`$ is an $`h`$-invariant set. So if $`g`$ is unbounded and $`\text{Fix}(g)\mathrm{}`$ then $`\text{Fix}(g)`$ is a single point fixed by $`h`$. If $`h`$ is also unbounded then $`\text{Fix}(g)=\text{Fix}(h)`$. Hence the following lemma implies that $`Gx_0=x_0`$ and again finishes the proof of Theorem B.
###### Lemma 5.8.
For any two non-elliptic semisimple elements $`g,hG`$ there is a chain $`g=g_1,g_2,\mathrm{},g_n=h`$ of non-elliptic semisimple elements, each commuting with its successor. Furthermore, the set of non-elliptic semisimple elements generates $`G`$.
###### Proof.
The first claim is equivalent to the connectivity of the Tits boundary, which is equivalent to the assumption on the rank. Both claims are well-known. ∎
###### Remark 5.9.
Observe that the somewhat restrictive assumption that $`B`$ is an $`L^p`$-related space is used only in the proof of Claim 5.5, the rest of the argument being in the context of general ucus Banach spaces.
## 6. Minimal Sets
Let $`B`$ be a strictly convex reflexive Banach space and $`G`$ a group acting on $`B`$ by affine isometries. Consider the ordered category $`𝒞`$ of non-empty closed convex $`G`$-invariant subsets of $`B`$ endowed with $`G`$-equivariant isometric maps and inclusion order. The goal of this section is to study minimal elements of $`𝒞`$ (regardless of whether they exist). In Section 7 we shall prove their existence, under certain conditions (see Corollary 7.5).
The Mazur–Ulam theorem states that a surjective isometry between (real) Banach spaces is affine. It is not known (and probably not true under no further assumptions) whether the analogous of the Mazur–Ulam theorem holds in the general context of convex subsets of Banach spaces. However, for subsets of strictly convex spaces it is obviously true:
###### Lemma 6.1.
Let $`CB`$ be a convex subset. Then every isometric map $`CB`$ is affine.
###### Proof.
It is enough to show that for all $`x,yC`$ and every $`0<t<1`$ the point $`p=tx+(1t)y`$ is determined metrically. This is true since by strict convexity
$$\overline{\mathrm{B}}(x,(1t)xy)\overline{\mathrm{B}}(y,txy)=\{p\}.$$
In particular the morphisms of $`𝒞`$ are affine. Another useful geometric property of closed convex sets in $`B`$ is the existence of a nearest point projection.
###### Lemma 6.2.
Let $`C`$ be a non-empty closed convex subset of $`B`$. Then for every $`xB`$ there exist a unique point $`\pi _C(x)C`$ such that $`x\pi _C(x)=d(x,C)`$.
###### Proof.
The uniqueness follows from strict convexity. By the Hahn-Banach theorem $`C`$ is weakly closed since it is closed and convex; therefore, by reflexivity and the Banach–Alaoğlu theorem we have a nested family $`C\overline{\mathrm{B}}(x,d)`$ of weakly compact sets as $`dd(x,C)`$; its intersection yields existence. ∎
The map $`\pi _C:BC`$ is called the *nearest point projection* on $`C`$. We remark that it is not continuous in general. It is continuous for uniformly convex Banach spaces and non-expanding for Hilbert spaces. Still, the distance between a point and its projection is always a $`1`$-Lipschitz function:
###### Lemma 6.3.
Let $`C`$ be a non-empty closed convex subset of $`B`$. Then the function $`x\pi _C(x)x`$ from $`B`$ to $`𝐑`$ is $`1`$-Lipschitz.
###### Proof.
For any $`x,yB`$
$$\pi _C(x)x\pi _C(y)x\pi _C(y)y+yx.$$
###### Lemma 6.4.
If $`C𝒞`$ is a minimal element, then any convex $`G`$-invariant continuous (or lower semi-continuous) function $`\phi :C𝐑`$ is constant.
###### Proof.
If $`\phi `$ were to assume two distinct values $`s<t`$, then $`\phi ^1((\mathrm{},s])`$ would be a strictly smaller element of $`𝒞`$. ∎
###### Lemma 6.5.
Let $`C,C^{}𝒞`$ with $`C`$ minimal. Then the nearest point projection $`\pi =\pi _C^{}|_C:CC^{}`$ is affine.
###### Proof.
For every $`x,yC`$ and $`t[0,1]`$, the definition of $`\pi `$ implies
(6.i)
$$\begin{array}{c}\pi \left(tx+(1t)y\right)\left(tx+(1t)y\right)\hfill \\ \hfill \left(t\pi (x)+(1t)\pi (y)\right)\left(tx+(1t)y\right)\\ \hfill t\pi (x)x+(1t)\pi (y)y.\end{array}$$
It follows that the function $`C𝐑`$, $`x\pi (x)x`$ is convex. Clearly it is $`G`$-invariant, and by Lemma 6.3 it is continuous, hence Lemma 6.4 implies that $`\pi (x)x`$ is constant on $`C`$. This constant must be $`d(C,C^{})`$; as both the right-hand side and the left-hand side in (6.i) equal $`d(C,C^{})`$, it follows that
$$\left(t\pi (x)+(1t)\pi (y)\right)\left(tx+(1t)y\right)=\pi \left(tx+(1t)y\right)\left(tx+(1t)y\right).$$
Therefore, by the uniqueness part of Lemma 6.2, $`t\pi (x)+(1t)\pi (y)`$ must be $`\pi \left(tx+(1t)y\right)`$. ∎
###### Lemma 6.6.
If $`C𝒞`$ is minimal and $`T:CB`$ is a $`G`$-equivariant affine map, then there exist a $`\varrho (G)`$-invariant vector $`bB`$ such that $`T(c)=c+b`$ for all $`cC`$.
###### Proof.
The map $`C𝐑`$, $`xTxx`$ is $`G`$-invariant, continuous and convex, hence by Lemma 6.4 it has a constant value $`d0`$. Since $`B`$ is strictly convex and $`C`$ is convex, the affine map $`\sigma (x)=Txx`$ from $`C`$ to the sphere of radius $`d`$ in $`B`$ must be constant. Its value $`b=\sigma (C)`$ is the desired ($`\varrho (G)`$-invariant) translation vector. ∎
###### Corollary 6.7.
The map $`\pi _C:CC^{}`$ from Lemma 6.5 is in fact a translation.
###### Corollary 6.8.
If $`C,C^{}𝒞`$ are minimal, then they are equivariantly isometric. Moreover, any equivariant isometry $`CC^{}`$ is a translation by a $`\varrho (G)`$-invariant vector.
###### Proof.
By Corollary 6.7, $`\pi _C^{}|_C:CC^{}`$ is an isometry; it is $`G`$-equivariant and hence onto by minimality of $`C^{}`$. The second claim follows from Lemma 6.1 and Lemma 6.6. ∎
## 7. Actions of Product Groups and Splitting
### 7.a.
The main goal of this section is to prove Theorem C. By Proposition 2.13 we may assume the affine action to be isometric with respect to a ucus norm on a Banach space $`B`$. The main step is the following theorem.
###### Theorem 7.1.
Let $`G=G_1\times G_2`$ be a product of topological groups with a continuous action by affine isometries on a uniformly convex Banach space $`B`$ without $`G`$-fixed point. Assume that the associated linear $`G`$-representation $`\varrho `$ does not almost have non-zero invariant vectors. Then there exists a non-zero $`\varrho (G_i)`$-invariant vector for some $`i\{1,2\}`$.
The proof of Theorem 7.1 uses minimal sets (in analogy to \[Mo2\]); notice that we are in the setting of Section 6 since uniformly convex spaces are reflexive and strictly convex \[BL, App. A\]. More precisely, we show:
###### Proposition 7.2.
Let $`G`$ and $`B`$ be as above. Then there exists a minimal non-empty closed convex $`G_1`$-invariant subset in $`B`$. In fact, any non-empty closed convex $`G_1`$-invariant subset contains a minimal such subset.
###### Proof of Theorem 7.1.
Proposition 7.2 provides a minimal non-empty closed convex $`G_1`$-invariant set $`CB`$. If there is no non-zero $`\varrho (G_1)`$-invariant vector, Lemma 6.6 (applied to $`G_1`$) shows that $`G_2`$ fixes every point of $`C`$. Since $`G_1`$ preserves $`C`$ and $`G`$ has no fixed point, $`C`$ cannot consist of a single point. Picking two distinct points $`x,yC`$ yields the non-zero $`\varrho (G_2)`$-invariant vector $`xy`$. ∎
Recall that uniform convexity is characterized by the positivity of the convexity modulus $`\delta `$ defined in Section 2.a. Moreover, $`\delta `$ is a positive, non-decreasing function which tends to zero at zero. Defining
$$\delta ^1(t)=sup\{\epsilon :\delta (\epsilon )t\},$$
$`\delta ^1`$ is easily seen to share the same properties. Furthermore, for every $`\epsilon >0`$, $`\delta ^1\delta (\epsilon )\epsilon `$.
###### Proof of Proposition 7.2.
Let $`C_0B`$ be any non-empty closed convex $`G_1`$-invariant subset; we will show that $`C_0`$ contains a minimal subset (if no initial $`C_0`$ was prescribed, one may choose $`C_0=B`$).
Pick any $`pC_0`$ and let $`C_1C_0`$ be the closed convex hull of the $`G_1`$-orbit of $`p`$. By Hausdorff’s maximal principle, we can chose a maximal chain $`𝒟`$ of non-empty closed convex $`G_1`$-invariant subsets of $`C_1`$. If $`b_C:=\pi _C(0)`$ is bounded as $`C`$ ranges over $`𝒟`$, then for some $`R>0`$ we have a nested family of non-empty sets $`\overline{\mathrm{B}}(0,R)C`$ which are weakly compact by reflexivity, Hahn–Banach theorem and Banach–Alaoğlu theorem. In particular the intersection $`𝒟`$ is non-empty, thus providing a minimal set for $`G_1`$. Therefore, we may from now on assume for a contradiction that the (non-decreasing) net $`R_C:=b_C`$ is unbounded over $`C𝒟`$. Let $`𝒟^{}𝒟`$ be the cofinal segment defined by $`R_C>0`$. We will obtain a contradiction by showing that for every compact $`KG`$, $`\mathrm{diam}(\varrho (K)\widehat{b}_C)`$ tends to zero along $`C𝒟^{}`$, where $`\widehat{b}_C=\frac{b_C}{R_C}`$.
Indeed, choose $`K_iG_i`$ compact with $`KK_1\times K_2`$ and let $`L=\mathrm{max}_{gK_1\times K_2}g0`$. The choice of $`b_C`$ implies $`gb_C0`$ and $`R_C\frac{b_C+gb_C}{2}`$ for all $`gG`$. Therefore, setting $`x=\frac{b_C}{gb_C}`$, $`y=\frac{gb_C}{gb_C}`$, the convexity modulus $`\delta _{C,g}:=\delta (xy)`$ gives
$$\begin{array}{c}R_C\frac{1}{2}b_C+\frac{1}{2}gb_C\frac{1}{2}(x+y)gb_C(1\delta _{C,g})(gb_Cg0+g0)\hfill \\ \hfill (1\delta _{C,g})(R_C+L)R_C(1+\frac{L}{R_C}\delta _{C,g})gK_1.\end{array}$$
Therefore $`\delta _C:=sup_{gK_1}\delta _{C,g}\frac{L}{R_C}0`$ along $`C𝒟^{}`$ and hence
$$\underset{gK_1}{sup}\frac{gb_Cb_C}{gb_C}\delta ^1(\delta _C)0.$$
Using $`gb_Cgb_Cg0+LR_C+L`$, it follows that
(7.i)
$$\underset{gK_1}{sup}\frac{gb_Cb_C}{R_C}0\text{along}C𝒟^{}.$$
On the other hand, for every $`gG_2`$, the function $`zgzz`$ is continuous, convex and $`G_1`$-invariant;therefore, it is bounded by $`gpp`$ on $`C_1`$. Setting $`L^{}=\mathrm{max}_{gK_2}gpp`$, it follows now that for all $`k=(g_1,g_2)K`$ we have
$$\begin{array}{c}R_C\varrho (k)\widehat{b}_C\widehat{b}_C=kb_Cb_Ck0g_1b_Cb_C+g_2g_1b_Cg_1b_C+L\hfill \\ \hfill g_1b_Cb_C+L^{}+L.\end{array}$$
Thus, in view of (7.i), $`\mathrm{diam}(\varrho (K)\widehat{b}_C)`$ goes to zero as claimed. ∎
###### Proof of Theorem C.
We adopt the notation and assumptions of that theorem; let $`\varrho `$ be the linear part of the action. Assume first $`n=2`$. Since we have in particular $`B^{\varrho (G)}=0`$, Corollary 2.9 yields a canonical splitting $`B=B^{\varrho (G_1)}B^{\varrho (G_2)}B_0`$ invariant under $`\varrho (G)`$. Decomposing the cocycle $`GB`$ along this splitting shows that up to affine isometry we may assume that the affine $`G`$-space $`B`$ splits likewise as affine product of affine spaces with corresponding linear parts. However, Theorem 7.1 shows that the resulting affine $`G`$-action on $`B_0`$ must have a fixed point since $`B_0^{\varrho (G_i)}=0`$. Therefore we obtain a $`G`$-invariant affine subspace $`G`$-isometric to $`B^{\varrho (G_1)}B^{\varrho (G_2)}`$ in $`B`$, as claimed.
In order to obtain the general case $`n2`$, we only need to observe that Corollary 2.9 applied to the product $`G_1\times _{i2}G_i`$ allows us to apply induction on $`n`$. ∎
###### Remark 7.3.
The above proof characterizes as follows the subspaces $`B_iB`$ appearing in the statement of Theorem C: Upon possibly replacing the $`B_i`$ with the corresponding linear subspace (which corresponds to replacing the cocycles with cohomologous cocycles), we have $`B_i=B^{\varrho (G_i^{})}`$ for $`G_i^{}=_{ji}G_j`$.
### 7.b. A More Geometric Approach to Theorem B and a Step Towards Conjecture 1.6
Before going on towards the superrigidity theorem, let us explain a more geometric, and seemingly more general, approach to prove $`(T_B)(F_B)`$, which is based on minimal sets. First we shall formulate a very general statement in the vein of Conjecture 1.6:
###### Theorem 7.4.
Let $`B`$ be a ucus Banach space and $`G`$ a topological group with property $`(T_B)`$ and compact Abelianization. Then for any continuous affine isometric action of $`G`$ on $`B`$ there is a minimal non-empty closed convex subset $`CB`$. Moreover $`\text{Aut}_G(C)`$ is trivial, $`CB^{}`$ and $`C`$ is unique up to translations by a $`\varrho (G)`$-invariant vector.
The proof of Theorem 7.4 relies on the following consequence of our discussion of minimal sets:
###### Corollary 7.5.
Let $`G`$ be a topological group with a continuous action by affine isometries on a uniformly convex Banach space $`B`$. Assume that the associated linear representation does not almost have non-zero invariant vectors. Then there exists a unique minimal non-empty closed convex $`G`$-invariant subset $`C_0B`$. Moreover, there are no non-trivial $`G`$-equivariant isometries of $`C_0`$.
###### Remark 7.6.
In view of the additional statement of Proposition 7.2, the set $`C_0`$ is contained in every non-empty closed convex $`G`$-invariant subset. Thus it is indeed the (non-empty) intersection of all those subsets.
###### Proof of Corollary 7.5.
For the existence of $`C_0`$, we may apply Proposition 7.2 if $`G=G_1\times 1`$ has no fixed point, or otherwise take such a fixed point for $`C_0`$. Both uniqueness and the additional statement follow now from Corollary 6.8. ∎
###### Proof of Theorem 7.4.
Since $`G`$ has compact Abelianization, the $`\varrho (G)`$-invariant subspace $`B^{}`$ is in fact $`G`$-invariant as an affine space, as the projection of the cocycle to $`B^{\varrho (G)}`$ must be a homomorphism. It follows that every minimal non-empty closed convex $`G`$-invariant set is contained in some coset of $`B^{}`$. The existence and uniqueness of such subset $`C`$ inside $`B^{}`$ follows from Corollary 7.5. The fact that any two such sets are different by a $`\varrho (G)`$-invariant vector is a consequence of Corollary 6.8. ∎
Let us now describe an alternative proof for Theorem B. Let $`B`$ be an $`L^p`$-related Banach space as in Theorem B. We reduce to the case where $`G`$ is connected, simply-connected and almost-simple as in Section 5. Now $`G`$ either contains a copy of $`\mathrm{𝐒𝐋}_3(k)`$ or of $`\mathrm{𝐒𝐩}(4,k)`$ that, in each case, contains a semidirect product $`HU`$ with the strong relative property $`(T_B)`$ (see Lemma 5.4 and the proof of Claim 5.5 ). We decompose $`B=B^{\varrho (U)}B^{}`$ according to that $`U`$ action (note that by Howe-Moore $`B^{\varrho (U)}=B^{\varrho (G)}`$). Then $`B^{}`$ is invariant under the affine action of $`G`$ and $`H`$ does not almost has invariant vectors in $`B^{}`$. Hence, by Corollary 7.5 there is a unique minimal non-empty closed convex $`H`$-invariant subset $`CB^{}`$ and it has no non-trivial automorphisms which commute with the $`H`$-action. Since, by Claim 5.5, the centralizer of $`H`$ is non-compact it follows by Howe–Moore that $`|C|=1`$. As in Section 5, we can now finish the proof using Lemma 5.8.
## 8. Induction and Superrigidity
Let $`\mathrm{\Gamma }<G=G_1\times \mathrm{}\times G_n`$ be a lattice in a product of $`n2`$ locally compact groups. Under an irreducibility assumption, the splitting theorem (Theorem C) implies a superrigidity result for uniformly equicontinuous affine $`\mathrm{\Gamma }`$-actions on superreflexive spaces $`B`$. As before such an action can be viewed as an affine isometric $`\mathrm{\Gamma }`$-action on a ucus Banach space $`B`$. It therefore suffices to apply the splitting theorem to the *induced* $`G`$-action on an *induced space* $`L^p(G/\mathrm{\Gamma },B)`$ (compare \[Sh\] for the Hilbertian case).
The goal of this section is to address the various (mostly technical) issues that arise when carrying out this programme. We begin by preparing for a statement (Theorem 8.3 below) that will then imply a more general form of Theorem D.
### 8.a.
Let $`G`$ be a locally compact group and $`\mathrm{\Gamma }<G`$ a lattice. The induction procedure will work smoothly if $`\mathrm{\Gamma }`$ is uniform (i.e. cocompact); in order to treat some non-uniform cases, one introduces the following.
###### Definition 8.1 (\[M5, III.1.8\]).
The lattice $`\mathrm{\Gamma }`$ is *weakly cocompact* if the $`G`$-representation $`L_0^2(G/\mathrm{\Gamma })`$, i.e. the canonical complement of the trivial representation in $`L^2(G/\mathrm{\Gamma })`$, does not almost have non-zero invariant vectors.
One verifies that any cocompact lattice is weakly cocompact. If $`G`$ has property $`(T)`$, then all its lattices are weakly cocompact. This also holds if $`G`$ is any connected semisimple Lie group (\[Bk\], compare also \[M5, III.1.12\]). By Remark 4.3, this definition does not depend on considering $`L^2(G/\mathrm{\Gamma })`$ rather than $`L^p(G/\mathrm{\Gamma })`$ for some other $`1p<\mathrm{}`$.
###### Definition 8.2 (See \[Sh, 1.II\]).
Let $`p>0`$. The lattice $`\mathrm{\Gamma }`$ is *$`p`$-integrable* if either (i) it is uniform; or (ii) it is finitely generated and for some (or equivalently any) finite generating set $`S\mathrm{\Gamma }`$, there is a Borel fundamental domain $`𝒟G`$ (with null boundary) such that
$$_𝒟\chi (g^1h)_S^p𝑑h<\mathrm{}gG,$$
where $`_S`$ is the word-length associated to $`S`$ and $`\chi :G\mathrm{\Gamma }`$ is defined by $`\chi ^1(e)=𝒟`$, $`\chi (g\gamma ^1)=\gamma \chi (g)`$.
This formulation is a bit awkward so as to include all uniform lattices since (ii) would otherwise fail when $`G`$ is not compactly generated. Condition (ii) holds (with any $`p1`$) for all irreducible lattices in higher rank semisimple Lie/algebraic groups, see \[Sh, §2\]; it holds likewise for Rémy’s Kac–Moody lattices \[Ry\].
Finally, given a product structure $`G=G_1\times \mathrm{}\times G_n`$, we say that a lattice $`\mathrm{\Gamma }<G`$ is *irreducible* if its projection to each $`G_i`$ is dense.
###### Theorem 8.3.
Let $`\mathrm{\Gamma }`$ be an irreducible lattice in a locally compact $`\sigma `$-compact group $`G=G_1\times \mathrm{}\times G_n`$. Assume that $`\mathrm{\Gamma }`$ is weakly cocompact and $`p`$-integrable for some $`p>1`$. Let $`B`$ be a ucus Banach space with a $`\mathrm{\Gamma }`$-action by affine isometries.
If the associated linear $`\mathrm{\Gamma }`$-representation does not almost have invariant vectors, then there is a $`\mathrm{\Gamma }`$-closed complemented affine subspace of $`B`$ on which the $`\mathrm{\Gamma }`$-action is a sum of actions extending continuously to $`G`$ and factoring through $`GG_i`$. (Compare Remark 1.9.)
Theorem 8.3 indeed implies Theorem D in the wider generality of weakly cocompact $`p`$-integrable lattices, since Proposition 2.13 allows us to assume that the topological vector space of Theorem D is in fact a ucus Banach space with a $`\mathrm{\Gamma }`$-action by affine isometries.
A (simpler) application of the same techniques implies the following result:
###### Theorem 8.4.
Let $`\mathrm{\Gamma }`$ be an irreducible lattice in a locally compact $`\sigma `$-compact group $`G=G_1\times \mathrm{}\times G_n`$. Assume that $`\mathrm{\Gamma }`$ is weakly cocompact and $`p`$-integrable for some $`p>1`$.
Then any homomorphism $`\mathrm{\Gamma }𝐑`$ extends continuously to $`G`$.
This result was established by Shalom in the case of cocompact lattices \[Sh, 0.8\] (actually, his proof holds in the setting of square-integrable lattices). It is therefore unsurprising that our reults imply the generalisation stated in Theorem 8.4 above (see the end of this section).
### 8.b. Induction
Throughout this section, $`G`$ is a locally compact second countable group and $`\mathrm{\Gamma }<G`$ a lattice. In particular, the Haar measure induces a standard Lebesgue space structure on $`G/\mathrm{\Gamma }`$.
###### Remark 8.5.
Even though Theorem 8.3 and Theorem D was stated in the more general setting of $`\sigma `$-compact groups, it is indeed enough to treat the second countable case: one can reduce to the latter by a structural result of Kakutani–Kodaira \[KK\] (the details of the straightforward reduction are expounded at length in \[Mo2\]).
Let $`B`$ be any Banach space and $`1<p<\mathrm{}`$. We consider the Banach space $`E=L^p(G/\mathrm{\Gamma },B)`$ as in Section 2.e.
###### Lemma 8.6.
If $`B`$ is uniformly convex or ucus, then so is $`E`$.
###### Proof.
This follows from a result of Figiel and Pisier, see Theorem 1.e.9 point (i) in \[LT\], Volume II. ∎
Suppose now that $`B`$ is endowed with a linear isometric $`\mathrm{\Gamma }`$-representation $`\varrho `$. Then $`E`$ can be canonically isometrically identified
(8.i)
$$EL^{[p]}(G,B)^{\varrho (\mathrm{\Gamma })}$$
with the space of those Bochner-measurable $`\mathrm{\Gamma }`$-equivariant function classes $`f:GB`$ such that $`f_B:G/\mathrm{\Gamma }𝐑`$ is $`p`$-integrable (the latter condition is symbolized by the notation $`L^{[p]}`$). Here, we choose to interpret $`\mathrm{\Gamma }`$-equivariance as $`f(g\gamma )=\varrho (\gamma )^1f(g)`$. The isomorphism (8.i) can be e.g. realized by restricting equivariant maps to any Borel fundamental domain $`𝒟G`$ for $`\mathrm{\Gamma }`$ since $`𝒟G/\mathrm{\Gamma }`$ as Lebesgue spaces. This identification allows us to endow $`E`$ with a continuous linear isometric $`G`$-representation by left multiplication. This $`G`$-representation is called the *induced representation*. If we choose a fundamental domain $`𝒟G`$ and consider the corresponding map $`\chi `$ as in Definition 8.2, then this $`G`$-representation reads as follows for $`fE=L^p(G/\mathrm{\Gamma },B)`$:
(8.ii)
$$(hf)(g\mathrm{\Gamma })=\varrho (\chi (g)^1\chi (h^1g))f(h^1g\chi (h^1g)\mathrm{\Gamma })$$
(a good indication that the model (8.i) is more natural!).
###### Lemma 8.7.
Assume $`\mathrm{\Gamma }`$ weakly cocompact in $`G`$. If the linear $`\mathrm{\Gamma }`$-representation does not almost have invariant vectors, then the induced linear $`G`$-representation does not either.
###### Proof.
The proof given by Margulis in the unitary case \[M5, III.1.11\] holds without changes (recalling that we can apply weak cocompactness in the $`L^p`$ setting by Remark 4.3). ∎
Suppose now that $`B`$ is endowed with an isometric $`\mathrm{\Gamma }`$-action – not necessarily linear anymore. We want to endow $`E`$ with a continuous affine isometric $`G`$-action by identifying $`E`$ with a space of $`\mathrm{\Gamma }`$-equivariant function classes $`GB`$ as before, except that equivariance is now understood with respect to the affine $`\mathrm{\Gamma }`$-action. Formally, there is nothing to change to the special case of linear action considered above; the action is defined by left $`G`$-translation of equivariant maps, so that via the natural identification we get for $`fE=L^p(G/\mathrm{\Gamma },B)`$ the action
(8.iii)
$$(hf)(g\mathrm{\Gamma })=\chi (g)^1\chi (h^1g)f(h^1g\chi (h^1g)\mathrm{\Gamma })$$
in complete analogy with (8.ii). However, the $`L^p`$ integrability property might be lost. The condition (ii) of Definition 8.2 is a straightforward sufficient condition to retain integrability; cocompactness of $`\mathrm{\Gamma }`$ is also enough, because it ensures that one can choose $`𝒟`$ in such a way that for any compact $`CG`$ the set $`\{\eta \mathrm{\Gamma }:𝒟\eta C\mathrm{}\}`$ is finite \[B2, VII §2 Ex. 12\]. Compare \[Sh, §2\] (and \[Mo2, App. B\]).
In conclusion, we may always consider the continuous *induced (affine) isometric $`G`$-action* on $`E`$ when $`\mathrm{\Gamma }`$ is $`p`$-integrable.
By construction, the linear part of the induced affine action coincides with the induced linear $`G`$-representation on $`E`$ considered earlier. If we denote by $`b:\mathrm{\Gamma }B`$ the cocycle of the original affine $`\mathrm{\Gamma }`$-action, then comparing (8.ii) with (8.iii) shows that the cocycle $`\stackrel{~}{b}:GE`$ of the induced affine action is given by
(8.iv)
$$\stackrel{~}{b}(h)(g\mathrm{\Gamma })=b\left(\chi (g)^1\chi (h^1g)\right).$$
Moreover, the correspondence $`b\stackrel{~}{b}`$ induces a (topological) isomorphism $`H^1(\mathrm{\Gamma },B)H^1(G,E)`$.
At this point, we record the following.
###### Proposition 8.8.
Keep the notation of this section.
1. If $`\mathrm{\Gamma }`$ has property $`(F_B)`$ then so does $`G`$.
2. If $`G`$ has property $`(F_E)`$ and $`\mathrm{\Gamma }`$ is $`p`$-integrable, then $`\mathrm{\Gamma }`$ has property $`(F_B)`$.
###### Proof.
For (1), consider any continuous affine isometric $`G`$-action on $`B`$; then there is a $`\mathrm{\Gamma }`$-fixed point $`bB`$. The corresponding orbit map $`GB`$ descends to a continuous map $`G/\mathrm{\Gamma }B`$. The image of the normalized invariant measure on $`G/\mathrm{\Gamma }`$ in $`B`$ being preserved by $`G`$, it follows from Lemma 2.14 that there is a $`G`$-fixed point.
For (2), consider an affine isometric $`\mathrm{\Gamma }`$-action on $`B`$ and endow $`E`$ with the induced affine action as in the discussion above. Then there is a $`G`$-fixed point $`fE`$. It follows from the description of $`E`$ as space of equivariant maps that $`f`$ is essentially constant and that its essential value is a $`\mathrm{\Gamma }`$-fixed point of $`B`$. ∎
### 8.c. Superrigidity
In order to prove Theorem 8.3, we now analyse the interplay between the induction constructions and the setting of irreducible lattices $`\mathrm{\Gamma }<G=G_1\times \mathrm{}\times G_n`$ as in the beginning of this Section 8. We will roughly imitate the arguments given by Shalom in \[Sh\] when he deduces Corollary 4.2 *ibid*.
Keep all the above notations and write $`G_i^{}=_{ji}G_j`$. First we observe that the irreducibility of $`\mathrm{\Gamma }`$ implies that for each $`i`$ it is a well-posed definition to consider the maximal (possibly zero) linear subspace $`B_iB`$ on which the linear $`\mathrm{\Gamma }`$-representation $`\varrho `$ extends to a continuous $`G`$-representation $`\varrho _i:GG_i𝐎(B_i)`$ factoring through $`G_i`$; moreover $`B_i`$ is automatically closed by maximality.
The induced space $`E`$ is ucus by Lemma 8.6. The isometric (affine) $`G`$-action on $`E`$ has no fixed point by the very same argument given to prove Proposition 8.8(2). On the other hand, the linear part does not have almost invariant vectors by Lemma 8.7. Thus Theorem C applies: There is a $`G`$-invariant closed complemented affine subspace $`\underset{¯}{E}E`$ and an affine isometric $`G`$-equivariant isomorphism $`\underset{¯}{E}E_1\mathrm{}E_n`$, where each $`E_i`$ is a ucus space with an affine isometric $`G`$-action factoring through $`GG_i`$. In view of Remark 7.3, there is no loss of generality in assuming that $`E_i`$ is the space of $`G_i^{}`$-fixed under the induced linear representation. One verifies readily the following:
###### Lemma 8.9.
The map $`B_iEL^{[p]}(G,B)^{\varrho (\mathrm{\Gamma })}`$ that to $`vB_i`$ associates the function $`GB`$ defined by $`g\varrho _i(g^1)v`$ yields an isometric isomorphism of (linear) $`G`$-spaces $`B_iE_i`$. ∎
Indeed, since the image of $`\mathrm{\Gamma }`$ in $`G_i`$ is dense, the Fubini–Lebesgue theorem implies that any map $`f:GB`$ in $`E`$ that is $`G_i^{}`$-invariant in the linear representation on $`E`$ is an orbit map as in the lemma.
At this point we observe that if the subspaces $`B_i`$ had trivial intersection, we would indeed have found a subspace $`B_iE_i`$ of $`B`$ on which the affine $`\mathrm{\Gamma }`$-action extends continuously to $`G`$ as requested. In general, we have a $`\mathrm{\Gamma }`$-equivariant affine map
$$E_iB_iB$$
induced by the maps of Lemma 8.9. Alternatively, we can think of this map as follows: The cocycle induced as in (8.iv) decomposes as a sum of cocycles $`\stackrel{~}{b}=\stackrel{~}{b}_i:GE`$, $`\stackrel{~}{b}_i:GG_iE_i`$, and in turn by Lemma 8.9 each $`\stackrel{~}{b}_i`$ is the cocycle induced under the correspondence (8.iv) from a cocycle $`b_i:\mathrm{\Gamma }B_i`$; the affine $`\mathrm{\Gamma }`$-action on $`B_i`$ is determined by the cocycle $`b_i`$. This completes the proof of Theorem 8.3. ∎
###### Remark 8.10.
As mentionned in Remark 1.9, the obstruction to extending the affine $`\mathrm{\Gamma }`$-action on some subspace of $`B`$ is confined within a compact group. Indeed, the only reason we might end up with a *sum* of action extending to $`G`$ through various $`G_i`$ rather than with a direct sum (which then extends globally to $`G`$) is the possibility that $`B_iB_j0`$ for some $`ij`$. But then the linear representation of $`\mathrm{\Gamma }`$ on $`B_iB_j`$ extends continuously to $`G`$ in two different ways, *both* through $`G_i`$ and through $`G_j`$. This may indeed happen but forces the image of $`\mathrm{\Gamma }`$ in $`𝐎(B_iB_j)`$ to be compact, see examples and discussion in \[Mo2\].
Let us only mention the most basic example: $`\mathrm{\Gamma }<G=G_1\times G_2`$ with $`G_i=𝐙\{\pm 1\}`$ and $`\mathrm{\Gamma }=𝐙^2\{\pm 1\}`$. Then $`\mathrm{\Gamma }`$ acts affinely isometrically without fixed point on $`B=𝐑`$ (by $`(n,m;\epsilon ).x=\epsilon x+n+m`$) and the associated linear representation does not almost have invariant vectors. However, it is easy to check that this action does not extend to $`G`$. Instead, it is a sum of actions extending to $`G_i`$ with sum map $`𝐑𝐑B=𝐑`$. Here $`B_1=B_2=B`$.
###### Proof of Theorem 8.4.
Recall that the space of homomorphisms $`\mathrm{\Gamma }𝐑`$ is precisely the space of affine isometric $`\mathrm{\Gamma }`$-actions on $`𝐑`$ with the trivial representation as linear part. By Remark 4.3, the $`G`$-representation on $`L_0^p(G/\mathrm{\Gamma })`$ does not almost have invariant vectors. Therefore, using $`p`$-integrable induction, one deduces Theorem 8.4 from Theorem C very exactly as Shalom deduced Theorem 0.8 in \[Sh\] from Theorem 3.1 in \[Sh\]. ∎
## 9. Appendix: Howe–Moore Theorem on Banach Spaces
In this appendix we sketch the proof of a version of the well known Howe–Moore theorem on vanishing of matrix coefficients for unitary representations, extended to the framework of ucus Banach spaces. This generalization is due to Yehuda Shalom (unpublished) and we state it here with a sketch of the proof for reader’s convenience.
###### Theorem 9.1.
Let $`I`$ be a finite set, $`k_i`$, $`iI`$ be local fields, $`𝐆_i`$ connected semisimple simply-connected $`k_i`$-groups, $`G_i=𝐆_i(k_i)`$ the locally compact group of $`k_i`$-points, and $`G=_{iI}G_i`$.
Let $`B`$ be a ucus Banach space and $`\varrho :G𝐎(B)`$ a continuous isometric linear representation, such that $`B^{\varrho (G_i)}=\{0\}`$ for each $`iI`$. Then all matrix coefficients $`c_{x,\lambda }(g)=\varrho (g)x,\lambda `$, $`xB`$, $`\lambda B^{}`$, vanish at infinity, i.e. $`c_{x,\lambda }C_0(G)`$.
Notice that we can (and will) assume that the $`𝐆_i`$ have no $`k_i`$-anisotropic factors, since the group of $`k_i`$-points of such factors are compact.
###### Proof of Theorem 9.1.
In a way of contradiction, assume that for some $`g_n\mathrm{}`$ in $`G`$, $`v\mathrm{S}(B)`$, $`\lambda \mathrm{S}(B^{})`$ one has
$$inf|\varrho (g_n)x,\lambda |=\epsilon >0.$$
We shall prove that at least one simple factor $`G_i`$ of $`G`$ has a non-trivial $`\varrho (G_i)`$ invariant vector.
Let $`G=KAK`$ be a Cartan decomposition of $`G`$ (here $`K=K_i`$ and $`A=A_i`$ where $`G_i=K_iA_iK_i`$ is the Cartan decomposition for $`G_i`$). We first show that without loss of generality one may assume $`g_nA`$.
###### Lemma 9.2 (KAK Reduction).
There exists a sequence $`a_n\mathrm{}`$ in the Cartan subgroup $`AG`$ and non-zero vectors $`y,zB`$ so that
$$\varrho (a_n)y\stackrel{w}{}z0.$$
where $`\stackrel{w}{}`$ denotes the weak convergence.
###### Proof.
Write $`g_n=k_na_nk_n^{}`$ where $`k_n,k_n^{}K`$ and $`a_nA`$. Then $`a_n\mathrm{}`$ because $`g_n\mathrm{}`$. Upon passing to a subsequence, $`k_n^{}k^{}K`$ and $`k_nkK`$. Denote
$$y_n=\varrho (k_n^{})x,y=\varrho (k^{})x,\mu _n=\varrho ^{}(k_n^1)\lambda ,\mu =\varrho ^{}(k^1)\lambda $$
where $`\varrho ^{}`$ is the dual (contragradient) $`G`$-representation on $`B^{}`$. Using the weak-compactness of the unit ball of $`B`$ we may also assume that
$$\varrho (a_n)y\stackrel{w}{}z.$$
We shall show that $`z,\mu =lim\varrho (g_n)x,\lambda `$ which is bounded away from zero, hence implying $`z0`$.
Recall that in a uc Banach space $`B`$ the weak and the strong topologies agree on the unit sphere $`\mathrm{S}(B)`$: indeed if $`y_n\stackrel{w}{}y`$ are unit vectors, then
$$1\delta (y_ny)y_n+y/2(y_n+y)/2,y^{}1.$$
Hence $`\delta (y_ny)0`$ and $`y_ny0`$. For the same reason we also have $`\mu _n\mu 0`$ in $`\mathrm{S}(B^{})`$. For an arbitrary $`\xi B^{}`$
$$|\varrho (a_n)y_n,\xi \varrho (a_n)y,\xi |y_ny\xi 0.$$
Hence $`\varrho (a_n)y_n\stackrel{w}{}z`$. In general, if $`z_n\stackrel{w}{}z`$ in $`B`$ and $`\mu _n\stackrel{w}{}\mu `$ in $`B^{}`$ then $`z_n,\mu _nz,\mu `$ because weakly convergent sequences are bounded in norm and
$`|z_n,\mu _nz,\mu |`$ $``$ $`|z_n,\mu _n\mu |+|z_nz,\mu |`$
$``$ $`(supz_n)\mu _n\mu ^{}+|z_nz,\mu |0.`$
Therefore
$$\varrho (g_n)x,\lambda =\varrho (a_nk_n^{})x,\varrho ^{}(k_n^1)\lambda =\varrho (a_n)y_n,\mu _nz,\mu $$
implying $`|z,\mu |\epsilon `$, which in particular means that $`z0`$. ∎
###### Lemma 9.3 (Generalized Mautner Lemma).
Suppose that $`\{a_n\}`$ and $`h`$ in $`G`$ satisfy $`a_n^1ha_n1_G`$ in $`G`$. If $`y,zB`$ are such that $`\varrho (a_n)y\stackrel{w}{}z`$ then $`\varrho (h)z=z`$. In particular, if $`\varrho (a_n)z=z`$ then $`\varrho (h)z=z`$.
###### Proof.
(Strong) continuity of $`\varrho `$ gives
$$\varrho (ha_n)y\varrho (a_n)y=\varrho (a_n^1ha_n)yy0$$
At the same time $`\varrho (a_n)y\stackrel{w}{}z`$ and $`\varrho (ha_n)y\stackrel{w}{}\varrho (h)z`$. Hence $`\varrho (h)z=z`$. ∎
We can now prove Theorem 9.1 in the case of $`G=\mathrm{𝐒𝐋}_2(k)`$ where $`k`$ is a local field. Assuming that $`\varrho :\mathrm{𝐒𝐋}_2(k)𝐎(B)`$ has some matrix coefficient not vanishing at infinity, we get by Lemma 9.2 a sequence $`a_n\mathrm{}`$ in $`A`$, and *non zero* vectors $`y,zB`$ with $`\varrho (a_n)y\stackrel{w}{}z`$.
Let $`H`$ be the unipotent (horocyclic) subgroup defined by $`H=\{hG:a_n^1ha_ne\}`$. It is normalized by $`a_n`$, and by Lemma 9.3 $`z`$ is a (non-trivial) $`\varrho (H)`$-invariant vector. We may assume that $`z=1`$. The matrix coefficient $`f(g)=\varrho (g)z,z^{}`$ is a continuous function on $`G`$, which is bi-$`H`$-invariant:
(9.i) $`f(gh)`$ $`=\varrho (g)\varrho (h)z,z^{}=\varrho (g)z,z^{}=f(g)`$
(9.ii) $`f(hg)`$ $`=\varrho (g)z,\varrho ^{}(h^1)z^{}=\varrho (g)z,z^{}=f(g)`$
for all $`gG`$ and $`hH`$. The proof can be now completed as in the original unitary Howe–Moore Theorem. By (9.i), $`f`$ can be viewed as a continuous function $`f_0`$ on the punctured plane $`G/H=k^2\{(0,0)\}`$, and by (9.ii), $`f_0`$ is constant on each horizontal line $`\mathrm{}_s=\{(t,s):tk\}`$, $`s0`$, where we identify $`H`$ with the upper triangular unipotent subgroup by choosing an appropriate basis for $`k^2`$. By continuity, $`f_0`$ is a constant on $`\{(t,0):t0\}`$. Since $`f_0(0,1)=f(e)=1`$ this constant is $`1`$.
This implies that $`z`$ is $`\varrho (A)`$-invariant because $`\varrho (a)z,z^{}=f(a)=f(e)=1`$ whilst $`z^{}`$ attains its norm only on $`z`$.
Thus $`z`$ is fixed by the upper triangular group $`AHG`$ and $`f`$ descends to a continuous function $`f_1`$ on the projective line $`𝐏(k^2)=G/AH`$. The $`H`$-action on $`𝐏(k^2)`$ has a dense orbit. Thus $`f_1`$ is constant $`1`$, and so is $`f`$:
$$\varrho (g)z,z^{}=f(g)=f(e)=1(gG)$$
Thus the unit vector $`z`$ is $`\varrho (G)`$-invariant, completing the proof in the case of $`G=\mathrm{𝐒𝐋}_2(k)`$.
The proof of the unitary Howe–Moore theorem for semisimple Lie group $`G=G_i`$ (c.f. Zimmer \[Z\], Margulis \[M5\]) relies only on the reduction to the Cartan subgroup (Lemma 9.2), the structure of such groups, the case of $`\mathrm{𝐒𝐋}_2(k)`$ and on Mautner Lemma. Thus the “unitary” argument can be applied almost *verbatim* to the present setup of ucus Banach spaces. ∎
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# Influence Of Collapsing Matter On The Enveloping Expanding Universe11footnote 1This work is dedicated to the memory of Professor Albert Einstein the greatest scientist of our time in commemoration of the postulation of the special theory of relativity 100 years ago. 22footnote 2PACS Codes:04.20.-q,11.10.Wx
## 1 Introduction
In several earlier papers Choudhury has constructed models where a Gidding-Strominger wormhole is placed at the center of an expanding universe and showed that, under certain conditions, the physical universe can expand with calculable acceleration. Following a different approach we now try to put a collapsing matter cloud at the center of the physical universe. We assume that the collapsing cloud generates a pressure which obeys the classical adiabatic gas law. This pressure penetrates the surrrounding universe. We add this pressure to the energy momentum tensor and compute the Hubble parameter and the deacceleration parameter, both of which are time dependent. The idea of the collapsing matter at the center of the universe is adopted from the treatment elaborated in the book of Weinberg . We have shown in section 1 how the $`R_c`$, the scale factor of the core, is calculated. The parametric solution depends on a quantity $`\psi `$. In section 2 we have we have calculated the pressure by using the classical adiabatic gas law. In section 3 we are assuming the enveloping universe to be homogeneous and isotropic, ignoring the core of the collapsing matter. We obtained the scale factor $`R(t)`$. We first calculated the Hubble parameter and the deacceleration parameter without incorporating the pressure from the core for large $`t_0`$ in section 4. We have shown that the enveloping universe is expanding with a constant deceleration. In section 5, we assumed that the pressure generated in the collapsing core is transmitted into the enveloping physical universe following classical Pascal’s law. We adhoc added this pressure to the energy momentum tensor. We then recalculated both the Hubble and the deacceleration parameter. Under certain restriction we have shown that the Hubble parameter can be positive and finite. The deacceleration parameter turns out to be fluctuating with time. In section 6 we discussed the significance of the result we obtained.
## 2 Pressure From Collapsing Matter
We assume that a certain amount of matter of the density $`\rho _c`$ at the core of the physical universe is collapsing in accordance with Weinberg’s treatment . For the moment we ignore the pressure. In the comoving coordinate system the invariant interval is given by
$$d\tau ^{c2}=dt^2+W^c(r,t)dr_{c}^{}{}_{}{}^{2}+V^c(r,t)(d\theta _c^2+sin^2\theta _cd\varphi _c^2).$$
(1)
With pressure negligible, we write
$$T^{c\mu \nu }=\rho ^cU^{c\mu }U^{c\nu },$$
(2)
with
$$U^{cr}=U^{c\theta }=U^{c\varphi }=0,U^{ct}=1.$$
(3)
In the above equations the index ”c” refers to collapsing entities. The Einstein’s field equation then becomes
$$R_{\mu \nu }^c=8\pi GS_{\mu \nu }^c.$$
(4)
with
$$S_{\mu \nu }^c=T_{\mu \nu }^c\frac{1}{2}g_{\mu \nu }^cT_{}^{c\lambda }{}_{\lambda }{}^{}.$$
(5)
Following Weinberg we can show that
$$W^c=R^{c2}(t)f^c(r)andV^c=R^{c2}(t)r^2.$$
(6)
with
$$f^c(r)=(1kr^2)^1.$$
(7)
The equation $`R^c(t)`$ satisfies
$$\dot{R}^{c2}(t)=\alpha [R^{c(1)}(t)1],$$
(8)
where
$$\alpha =\frac{8\pi G}{3}\rho ^c(0).$$
(9)
The solution is the parametric equation of cycloid
$$t=\frac{\psi +sin\psi }{2\sqrt{\alpha }}.$$
(10)
and
$$R^c=\frac{1}{2}(1+cos\psi ).$$
(11)
We now introduce a simplified approximation assuming $`\psi `$ to be a small quantity. Then t takes the form
$$t=\frac{\psi }{\sqrt{\alpha }},$$
(12)
and
$$R^c=cos^2(\frac{\sqrt{\alpha }t}{2}).$$
(13)
Before total collapse we assume that the core behaves as an adiabatic gas . The volume of the collapsing matter can be shown to be
$$\tau ^c=2\pi ^2R^{c3}(t).$$
(14)
This thermodynamical core satisfies the adiabatic gas law
$$P^c\tau ^{c\gamma }=constant=B_1.$$
(15)
This yields
$$P^c=Bcos^{6\gamma }(\frac{\sqrt{\alpha }t}{2}),$$
(16)
We conjecture that this pressure is transmitted into the surrounding expanding physical universe according to Pascal’s law.
## 3 Enveloping Space
We assume that the matter which surrounds the collapsing matter is homogeneous and isotropic. The extension of physical universe is so immense that if there is any violation of the properties of homogeneity and isotropy of the collapsing matter it can be ignored in our consideration. The space-time interval can be given by
$$d\tau ^2=dt^2+R^2(t)[f(r)dr^2+r^2d\theta ^2+sin^2\theta d\varphi ^2).$$
(17)
The Einstein equation for the enveloping space would be
$$R_{\mu \nu }=8\pi GS_{\mu \nu },$$
(18)
where
$$S_{\mu \nu }=T_{\mu \nu }\frac{1}{2}g_{\mu \nu }T_{}^{\lambda }{}_{\lambda }{}^{}.$$
(19)
with
$$T_{\mu \nu }=Pg_{\mu \nu }+(P+\rho )U_\mu U_\nu .$$
(20)
In Eq.(21) the U’s are given by
$$U^r=U^\theta =U^\varphi =0,U^t=1.$$
(21)
The only nonvanishing compopnents of the Ricci tensor are
$$R_{tt}=\frac{3\ddot{R}}{R},$$
(22)
$$R_{rr}=\frac{f^{}(r)}{rf(r)}+\ddot{R}Rf(r)+2\dot{R}^2f(r),$$
(23)
$$R_{\theta \theta }=1\frac{1}{f(r)}+\frac{rf^{}(r)}{f^2(r)}+\ddot{R}Rr^2+2\dot{R}^2r^2,$$
(24)
and
$$R_{\varphi \varphi }=R_{\theta \theta }sin^2(\theta ).$$
(25)
The righthand side of the components of Eq.(3.2) becomes
$$S_{tt}=\frac{\gamma }{2}(\rho +3P),$$
(26)
$$S_{rr}=\frac{\gamma }{2}R^2f(P+\rho ),$$
(27)
$$S_{\theta \theta }=\frac{\gamma }{2}R^2r^2(P+\rho ),$$
(28)
and
$$S_{\varphi \varphi }=\frac{\gamma }{2}R^2r^2sin^2\theta (P+\rho ).$$
(29)
In the above equations we have set
$$\gamma =8\pi G.$$
(30)
Taking the rr-component of the Einstein equation we get
$$\frac{f^{}(r)}{rf(r)}+\ddot{R}Rf(r)+2\dot{R}^2f(r)=\frac{\gamma }{2}R^2f(P+\rho ),$$
(31)
Here we make a conjecture that the pressure and density of the physical universe only depends on time t. We find that the Eq.(32) yields
$$2\sigma +\ddot{R}R+2\dot{R}^2=\frac{\gamma }{2}R^2(P+\rho ),$$
(32)
where we define a constant $`\sigma `$ by the relation
$$\frac{f^{}(r)}{rf^2(r)}=2\sigma .$$
(33)
From $`\theta \theta `$-component setting
$$2\beta =\frac{1}{r}^2\frac{1}{r^2f(r)}+\frac{f^{}(r)}{rf^2(r)},$$
(34)
we get
$$2\beta +\ddot{R}R+2\dot{R}^2=\frac{\gamma }{2}R^2(P+\rho ).$$
(35)
Subtracting from Eq.(33), Eq.(36) we obtain
$$2(\sigma \beta )=2R^2\gamma P.$$
(36)
The above equation yields
$$P=\frac{\lambda }{\gamma R^2},$$
(37)
If we conjecture that $`\sigma >\beta `$, P will always be positive.
## 4 Hubble and Deacceleration Parameters Without The Influence Of A Collapsing Core
In this section we start with the assumption that the collapsing core does exist, but the physical universe is expanding. To compute the Hubble parameter we turn to the tt-component of the Einstein Eq.(19). We get
$$\ddot{R}=\frac{\gamma }{6}R(\rho +3P).$$
(38)
Combinig the above equation with Eq.(3.11)we get
$$H(t)=\frac{\dot{R}}{R}=\sqrt{(\frac{\gamma \rho }{3}+2\gamma P\frac{2\sigma }{R^2})}.$$
(39)
Since the physical universe is, by assumption, expanding,we have for large $`tt_0`$, $`R0`$.
$$H(t_0)\sqrt{\frac{\gamma \rho }{3}+2\gamma P}.$$
(40)
For large $`t_0`$, the deacceleration parameter becomes:
$$q_0(t_0)=\frac{\ddot{R}R}{\dot{R}^2}=\frac{\frac{\gamma }{6}(\rho +3P)}{\gamma (\frac{\rho }{3}+2P)\frac{2\sigma }{R^2}}.$$
(41)
Since for large R, P vanishes, we obtain
$$q_0\frac{1}{2}.$$
(42)
Thus the universe is decelerating at a constant rate. But all cosmic observation indicates that at the present time the universe is expanding with acceleration . Instead of introducing the pressure from an expanding wormhole, we generate the pressure from a collapsing core as described in section 2.
## 5 Influence Of Pressure From Collapsing Matter
We now add the pressure which we conjectured in section 2, into the energy momentum tensor. This pressure originates from the core which is collapsing adiabatically. We take this pressure to be of the form $`P_c`$, derived in the Eq.(17). This pressure is transmitted into the enveloping universe by classical Pascal’s law. We add this amount to the pressure in Eq.(40) and get
$$H(t)=\frac{\dot{R}}{R}=\sqrt{\frac{\gamma \rho }{3}+2\gamma (P+P_c)\frac{2\sigma }{R^2}}$$
(43)
Now from Eq.(3.16) $`P0`$ as $`R\mathrm{}`$, we find for large time $`t_0`$
$$H_0(t_0)=\sqrt{\gamma (\frac{\rho }{3}+2Bcos^{6\gamma }(\frac{\sqrt{\alpha t_0}}{2}))}.$$
(44)
We assume that the quantity $`\sqrt{\alpha t_0}/2`$ is small even for large $`t_0`$ because $`\alpha `$ is very small. By choosing $`\gamma `$ appropriately we can make $`H_0(t_0)`$ fluctuate with time. There is a possibility that $`H_0`$ may become very large. In order to avoid such consequences we have to improve the solution of our approximation of Eqs.(11)and (12). With the addition of $`P_c`$ we find that the time dependent deacceleration parameter for an expanding universe becomes
$$q_0(t_0)=\frac{1}{2}\frac{\rho cos^{6\gamma }(\frac{\sqrt{\alpha }t_0}{2})+6B}{\rho cos^{6\gamma }(\frac{\sqrt{\alpha }t_0}{2})+12B}.$$
(45)
If we choose $`\gamma =1/6`$, we can write
$$q_0(t_0)=\frac{1}{2}[1\frac{6B}{\rho cos(\frac{\sqrt{\alpha }t_0}{2})+12B_1}].$$
(46)
This parameter is time dependent. The universe is expanding with acceleration if it satisfies the following condition:
$$6B>\rho cos(\frac{\sqrt{\alpha }t_0}{2})+12B$$
(47)
In all other times the universe decelerates.
## 6 Concluding Remarks
We have shown here that if we take the pressure and density of an expanding physical universe which are only time dependent, the Hubble time dependent parameter is positive. The deacceleration parameter for such a universe turns out to be positive, indicating that the universe decelerates. As an alternative idea we have put a gravitationally collapsing body at the center of the expanding universe and conjectured that the collapsing matter generates a pressure according to adiabatic gas law. This pressure transmits uniformly across the enveloping physical universe following Pascal’s law of fluid pressure. For large $`t_0`$ satisfying the condition $`\sqrt{\alpha }t_0`$ to be small, the Hubble parameter stays positive and finite indicating expansion of the physical universe. The deacceleration parameter fluctuates with time. Under certain condition the universe accelerates.
## 7 References
1. A. L. Choudhury, Hadronic J.,23, 581 (2000).
2. L. Choudhury and H. Pendharkar, Hadronic J., 24, 275 (2001).
3. A. L. Choudhury, Hadronic J.,27, 387 (2004).
4. A. L. Choudhury: Wormhole core, extra dimensions,and physical universe; arXiv:gt-qc/0405135v1,27 May 2004.
5. S. B. Giddings and A. Strominger, Nucl. Phys. B307, 854 (1988).
6. S. Weinberg: Gravitation and Cosmology, J. Wiley and Sons, Inc., 342 (1972).
7. N. Bahcall, J. P. Ostriker, S. Perlmutter, and P. J. Steinhardt, Science, 284, 1481 (1999).
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# Fractal upper bounds on the density of semiclassical resonances
## 1. Introduction
Let $`P=h^2\mathrm{\Delta }_g+V(x)`$ be a self-adjoint Schrödinger operator on a compact Riemannian $`n`$-manifold, $`(X,g)`$, $`V𝒞^{\mathrm{}}(X;)`$. The spectral asymptotics as $`h0`$ are given by the celebrated Weyl law – see and for recent advances and numerous references. If we assume that the zero energy surface is nondegenerate,
$$p\stackrel{\mathrm{def}}{=}|\xi |_g^2+V(x)=0dp0,$$
then
(1.1)
$$|\mathrm{Spec}(\mathrm{P})D(0,Ch)|=𝒪(h^{n+1}),$$
where $`D(0,r)=\{z:|z|<r\}`$, though of course in this case the eigenvalues are all real – see §6.1 for yet another proof of this well known result.
Let $`H_p`$ be the Hamilton vector field of $`p`$ on $`T^{}X`$, locally given by
$$H_p=\underset{j=1}{\overset{n}{}}\frac{p}{\xi _j}_{x_j}\frac{p}{x_j}_{\xi _j},(x,\xi )T^{}^n.$$
When the flow, $`\mathrm{exp}tH_p:p^1(E)p^1(E)`$, has the property that the set of its closed orbits has Liouville measure zero on $`p=0`$, then we have the infinitesimal version of the Weyl law:
(1.2)
$$|\mathrm{Spec}(\mathrm{P})D(0,Ch)|=\frac{2Ch}{(2\pi h)^n}_{p(x,\xi )=0}𝑑(x,\xi )+o(h^{n+1}),$$
where $`d`$ is the Liouville measure on $`p=0`$, that is $`ddp=dxd\xi `$. This result is the mathematical starting point of many recent investigations, mostly in physical literature, of the finer structure of the spectrum and its relation to classical dynamics – see and references given there.
When the manifold is non-compact the situation is dramatically different. The simplest case is that of a manifold which is Euclidean outside of a compact set and $`V+1𝒞_\mathrm{c}^{\mathrm{}}(X;)`$. The discrete eigenvalues of $`P`$ are replaced by quantum resonances which are defined as the poles of the meromorphic continuation of
$$(Pz)^1:𝒞_\mathrm{c}^{\mathrm{}}(X)𝒞^{\mathrm{}}(X),Imz>0,$$
and we denote the set of resonances by $`\mathrm{Res}(P(h))`$. The basic physical interpretation is that a resonance at $`z=E_0i\mathrm{\Gamma }/2`$ corresponds to a state with time evolution given by $`\mathrm{exp}(itE_0/h\mathrm{\Gamma }t/2h)`$, and to a Breit-Wigner peak in energy density, $`\mathrm{\Gamma }/((EE_0)^2+\mathrm{\Gamma }^2/4)`$. This intuitive picture becomes however complicated when many resonances are present which is natural for $`h`$ is small – see and ,, for recent results and references.
Here we provide upper bounds for the number of resonances of $`P`$ in $`D(0,Ch)`$. The main result (Theorem 3) states that for classical Hamiltonians $`p`$ with hyperbolic flow on $`p=0`$,
(1.3)
$$|\mathrm{Res}P(h)D(0,Ch)|=𝒪(h^\nu ),$$
where $`2\nu +1`$ is essentially the dimension of the trapped (non-wandering) set in $`p^1(0)`$,
$$K\stackrel{\mathrm{def}}{=}\{(x,\xi )T^{}X:p(x,\xi )=0,\mathrm{exp}(tH_p)(x,\xi )\to ̸\mathrm{},t\pm \mathrm{}\}.$$
In the case of a compact manifold $`\nu =n1`$ so that (1.3) reduces to (1.1). By dimension we always mean the Minkowski dimension
$$m_0=2n1sup\{d:\underset{ϵ0}{lim\; sup}ϵ^d\mathrm{vol}(\{\rho p^1(0):d(\rho ,K)<ϵ\})<\mathrm{}\}.$$
A simple example is provided by a three bump potential shown in Fig.1.
The first estimate of this form was proved by the first author in \[29, Theorems 4.6, 5.5, and 5.7\]: there exists constants $`C_0,C_1>0`$, such that for $`\delta _0>0`$ fixed and small enough
(1.6)
$$\begin{array}{c}|\mathrm{Res}(P(h))\{z:|z|<\delta ,Imz>\mu \}|C_1\delta \left(\frac{h}{\mu }\right)^n\mu ^{\frac{1}{2}\stackrel{~}{m}},\\ C_0h\mu 1/C_0,C_0h^{\frac{1}{2}}\delta \delta _0,0<h<1/C_0,\end{array}$$
where now $`\stackrel{~}{m}`$ is any number greater than the dimension of the trapped set in
$$p^1([2\delta _0,2\delta _0]).$$
In homogeneous situations, such as for instance obstacle scattering, $`\stackrel{~}{m}=m+1`$. When $`\mu =C_0h`$, the improvement in Theorem 3 lies in allowing $`\delta h`$, which is the natural limit for this type of spectral estimates.
Earlier, non-geometric, bounds on the number of resonances (scattering poles) were obtained by Melrose , and the second author ,. In the case of convex co-compact Schottky quotients (and any convex co-compact quotients in dimension two) the analogue of (1.3) was proved in using zeta function techniques, improving earlier estimates of the proof of which was largely based on . These technique gave similar results for the zeros of zeta functions of rational maps ,, in which case the dimension of the trapped set becomes essentially the dimension of the Julia set.
Numerical investigations in different settings of semiclassical three bump potentials , (see Fig.2), Schottky quotients , three disc scattering , and Cantor-like Julia sets for $`zz^2+c`$, $`c<2`$ , suggest that for $`\mu Ch`$ and $`\delta 1`$ the estimate (1.6) is optimal. A different model was recently considered in : quantum resonances were defined using an open quantum map with a classical “trapped set” corresponding to $`K`$ intersected with a hypersurface transversal to the flow – see Fig.5. The numerical results and a simple linear algebraic toy model suggest that the fine estimate (1.3) is optimal – see Fig.3. A similar model was also used in where the fractal Weyl law gave corrections to the applications of random matrix theory to open quantum systems.
We now state the general assumptions on the operator $`P`$. We reiterate that the simplest case to keep in mind is
$$P=h^2\mathrm{\Delta }+V(x)1,V𝒞_\mathrm{c}^{\mathrm{}}(^n).$$
In general we consider
$$P(h)\mathrm{\Psi }^{0,2}(X),P(h)=P(h)^{},$$
where the calculus of semiclassical pseudodifferential operators is reviewed in §3.1,
(1.10)
$$\begin{array}{c}P(h)=p^w(x,hD)+hp_1^w(x,hD;h),p_1S^{0,2}(T^{}X),\\ |\xi |Cp(x,\xi )\xi ^2/C,p=0dp0,\\ R,u𝒞^{\mathrm{}}(XB(0,R))P(h)u(x)=Q(h)u(x),Q(h)u=\underset{|\alpha |2}{}a_\alpha (x;h)(hD_x)^\alpha u,\end{array}$$
where $`a_\alpha (x;h)=a_\alpha (x)`$ is independent of $`h`$ for $`|\alpha |=2`$, $`a_\alpha (x;h)C_b^{\mathrm{}}(^n)`$ are uniformly bounded with respect to $`h`$, here $`C_b^{\mathrm{}}(^n)`$ denotes the space of $`C^{\mathrm{}}`$ functions on $`^n`$ with bounded derivatives of all orders,
(1.13)
$$\begin{array}{c}\underset{|\alpha |=2}{}a_\alpha (x)\xi ^\alpha (1/c)|\xi |^2,\xi ^n,\text{ for some constant }c>0\text{}\\ \text{ }_{|\alpha |2}a_\alpha (x;h)\xi ^\alpha \xi ^21\text{ uniformly with respect to }h\text{ as }|x|\mathrm{}.\end{array}$$
We also need the following analyticity assumption in a neighbourhood of infinity: there exist $`\theta [0,\pi ),`$ $`ϵ>0`$ and $`RR_0`$ such that the coefficients $`a_\alpha (x;h)`$ of $`Q(h)`$ extend holomorphically in $`x`$ to
$$\{r\omega :\omega ^n,\text{dist}(\omega ,𝐒^n)<ϵ,r,|r|>R,\text{arg}r[ϵ,\theta _0+ϵ)\},$$
with (1.13) valid also in this larger set of $`x`$’s. We remark that in , the operators were required to be globally real analytic but the conditions at infinity were much more general. A particularly nice feature of the theory developed in is allowing arbitrary homogeneous polynomials as potentials (see \[29, (c.31)-(c.33)\]).
The first theorem we present fits naturally in the methodology of this paper. It is a slight generalization of a result of Martinez which in turn was a $`𝒞^{\mathrm{}}`$ version of a similar result in , already implicit in . As discussed at the end of §4 it is essentially optimal.
###### Theorem 1.
Suppose that $`X`$ is non compact and euclidean outside of a compact set, and that the operator $`P\mathrm{\Psi }^{0,2}(X)`$ satisfies the assumptions below. Suppose in addition that no orbit of $`H_p`$ on $`p^1(0)`$ is trapped:
(1.14)
$$Kp^1(0),T_K,(x,\xi )K\mathrm{exp}(tH_p(x,\xi ))K,t>T_K.$$
Then, for any $`M>0`$ there exists $`h(M)`$ such that for $`0<h<h(M)`$,
(1.15)
$$|\mathrm{Res}(P)D(0,Mh\mathrm{log}(1/h))|=\mathrm{}.$$
Here $`\mathrm{Res}(P)`$ denotes the set of resonances of $`P`$ defined in some $`h`$-independent neighbourhood of $`0`$.
Before stating the main result which requires hyperbolicity of the flow on the energy surface we first give the general upper bound which generalizes slightlyHere we consider an arbitrary manifold in the compact part. It is clear that the methods of easily allow this type of generalization and the point is that our method is different and more robust. the results of Bony which in turn generalized earlier results of Petkov and the second author ,.
###### Theorem 2.
Suppose that $`X`$ is Euclidean outside a compact set, and that the operator $`P\mathrm{\Psi }^{2,0}(X)`$ satisfies the general assumptions below. Then, as $`h0`$,
(1.16)
$$|\mathrm{Res}(P)D(0,Ch)|=𝒪(h^{n+1}).$$
The estimate (1.16) is also optimal in the same way that the analogous estimate for eigenvalues of a self-adjoint operator with a compact and smooth energy surface. That follows for instance from applying \[24, Corollary, §5\].
The basic hyperbolicity assumption at an energy $`E`$ can be stated as follows: for $`\rho p^1(E)`$ lying in a neighbourhoood of the trapped set $`K_E`$ we have,
(1.17)
$$\begin{array}{cc}& T_\rho (p^1(E))=H_p(\rho )E_+(\rho )E_{}(\rho ),dimE_\pm (\rho )=n1,\hfill \\ & p^1(E)\rho E_\pm (\rho )T_\rho (p^1(E))\text{ is continuous,}\hfill \\ & d(\mathrm{exp}tH_p)_\rho (E_\pm (\rho ))=E_\pm (\mathrm{exp}tH_p(\rho )),\hfill \\ & \lambda >0d(\mathrm{exp}tH_p)_\rho (X)Ce^{\pm \lambda t}X,\text{ for all }XE_\pm (\rho )\text{}t0\text{.}\hfill \end{array}$$
An example of a potential satisfying this assumption at a range of non-zero energies is given in Fig.1 – see and \[29, Appendix c\]. Following we will formulate a weaker dynamical hypothesis in §7.
The main result of this paper is
###### Theorem 3.
Suppose that $`P(h)`$ satisfies our general assumptions (1.10), the flow of $`H_p`$ near zero energy is hyperbolic in the sense of (1.17), or the weaker sense given in §7.1, and that the trapped set at zero energy has Minkowski dimension $`m_0=2\nu _0+1`$. Then for any $`\nu >\nu _0`$, and $`C_0>0`$ there exists $`C_1`$ such that
(1.18)
$$|\mathrm{Res}(P(h))D(0,C_0h)|C_1h^\nu .$$
When the trapped is set is of pure dimension, $`\nu `$ can be replaced by $`\nu _0`$.
A sharp rigorous lower bound is known when $`K`$ is an isolated hyperbolic trajectory. A very precise asymptotic description of resonances in that case is given in and it implies (1.18) with $`\nu =0`$. In spite of the convincing numerical evidence cited above no rigorous examples with non-integral values $`\nu _0`$ are known. A recent indication of the delicate nature of lower bounds for resonances was given in where a class of complex compactly supported potentials in $`^3`$ with no resonances at all. We have no reasonable hope of obtaining any analogue of (1.2) at the present moment.
The methods of this paper apply also to a simpler problem of operators with complex absorbing barriers. Let $`V𝒞_\mathrm{c}^{\mathrm{}}(B(0,R_0);)`$ be a potential for which $`H_p`$, $`p=\xi ^2+V(x)1`$, has hyperbolic flow on $`p=0`$, for instance a “three bump” potential . Now let $`W𝒞^{\mathrm{}}(^n)`$ satisfy
$$W(x)0,\mathrm{supp}WB(0,R_0)=\mathrm{},W(x)2C_0>0\text{ for }|x|>2R_0\text{.}$$
Consider then (see and references given there)
$$\stackrel{~}{P}(h)=h^2\mathrm{\Delta }+V(x)iW(x).$$
The spectrum of this non-selfadjoint operator lies in $`\overline{}_{}`$ and we have the exact analogue of (1.3):
$$|\mathrm{Spec}\stackrel{~}{P}(h)D(0,Ch)|=𝒪(h^\nu ).$$
Acknowledgements. The second author would like to thank Jean-Yves Chemin and Jean-Michel Bony for useful discussions, the National Science Foundation for partial support under the grant DMS-0200732, and École Polytechnique for its generous hospitality in Fall 2004.
## 2. Outline of the proof
To prove the main result on fractal upper bounds (Theorem 3) we first develop methods for proving the natural results on the absence of resonances (Theorem 1) and on general upper bounds at non-degenerate energies (Theorem 2). In this section we present the general ideas. All of them have origins in other works and pointers to the literature will be given in corresponding sections.
The absence of resonances for operators with $`𝒞^{\mathrm{}}`$ coefficients in domains of size $`h\mathrm{log}(1/h)`$ around an energy level hold under a nontrapping condition:
$$ϵ_0>0,Kp^1(0),T_K,(x,\xi )K\mathrm{exp}(tH_p(x,\xi ))K,t>T_K.$$
This implies the existence of an escape function in a neighbourhood of $`p^1(0)`$:
$$G_1𝒞^{\mathrm{}}(T^{}X),H_pG_1(x,\xi )c_0>0,\text{for }|p(x,\xi )|<ϵ_0\text{.}$$
The resonances of $`P`$ are given by the eigenvalues of the deformed operator $`P_\theta `$. In the case of $`P=h^2\mathrm{\Delta }+V(x)`$ with $`V`$ analytic in a conic neighbourhood of $`^n`$,
$$V(x)+10,|x|\mathrm{},$$
the scaled operator is simply
$$P_\theta =h^2e^{2i\theta }\mathrm{\Delta }+V(e^{i\theta }x),$$
and it behaves as $`h^2e^{2i\theta }\mathrm{\Delta }1`$ near infinity. For $`\theta >0`$ that last operator is clearly invertible.
We can introduce a modified $`G=\chi G_1`$, $`\chi 𝒞_\mathrm{c}^{\mathrm{}}(X)`$, so that for
$$\theta ϵMh\mathrm{log}(1/h),$$
we have
$$|Rep_\theta |<\delta Imp_\theta +ϵH_pGc_0ϵ,p_\theta =\sigma (P_\theta ).$$
The operators $`\mathrm{exp}(\pm ϵG^w(x,hD)/h)`$ are now pseudodifferential operators in a mildly exotic class and we consider
$$P_{\theta ,ϵ}=e^{ϵG^w/h}P_\theta e^{ϵG^w/h}.$$
The spectrum of $`P_\theta `$ in $`D(0,Mh\mathrm{log}(1/h))`$ is the same as that of $`P_{\theta ,ϵ}`$ but the properties of $`G`$ imply that
$$P_{\theta ,ϵ}^1C/ϵ$$
showing that in fact there is no spectrum in $`D(0,M^{}h\mathrm{log}(1/h))`$. This approach allows us to obtain the absence of resonances very directly.
In Theorem 2 we show that if $`0`$ is a non-critical energy level then
$$|\mathrm{Res}PD(0,Ch)|=𝒪(h^{n+1}).$$
The proof follows from a “robust” proof of the same estimate for an operator with a compact resolvent (for instance, an elliptic operator on a compact manifold). Let $`P`$ be such an operator, say, $`P=h^2\mathrm{\Delta }_g1`$, on a compact Riemannian manifold. We would like to consider a modified operator
$$\stackrel{~}{P}(h)\stackrel{\mathrm{def}}{=}P(h)iMh\psi (MP(h)/h),\psi 𝒞_\mathrm{c}^{\mathrm{}}(),$$
whose “symbol”, $`piMh\psi (Mp/h)`$, has the absolute value bounded from below by $`Mh/2`$ everywhere. That does not make sense at first since
$$\psi (MP(h)/h)$$
is not an $`h`$-pseudodifferential operator. To remedy this we construct a second microlocal calculus with a new Planck constant $`\stackrel{~}{h}1/M`$. The new operator $`\stackrel{~}{P}(h)`$ becomes elliptic in this calculus and for $`\stackrel{~}{h}`$ small enough, independent of $`h`$, it is invertible. We then have
$$(P(h)z)^1=(I+K(z))^1(\stackrel{~}{P}(h)z)^1,K(z)\stackrel{\mathrm{def}}{=}i(\stackrel{~}{P}(h)z)^1Mh\psi (P(h)/h)$$
and the eigenvalues of $`P(h)`$ near $`0`$ coincide with the zeros of $`det(I+K(z))`$. The zeros of this determinant are the same as the zeros of a determinant $`det(I+R(z))`$ where $`R(z)`$ is a finite rank operator with the rank proportional to $`h^n`$ times the volume of the support of $`\psi (Mp/h)`$. That gives estimates on the deteminant which imply (1.1). A slight modification of this argument is needed to obtain Theorem 2.
We now assume that the flow of $`H_p`$ is hyperbolic and introduce the sets
(2.1)
$$\mathrm{\Gamma }_\pm \stackrel{\mathrm{def}}{=}\{(x,\xi )T^{}X:p(x,\xi )=0,\mathrm{exp}(tH_p)(x,\xi )\to ̸\mathrm{},t\mathrm{}\},$$
depicted in a simple case in Fig.4. The trapped set at zero energy is
(2.2)
$$K=\mathrm{\Gamma }_+\mathrm{\Gamma }_{}.$$
If we assume that $`K\overline{\mathrm{\Gamma }_\pm K}`$, that is $`K`$ has no component isolated from infinity, then $`K`$ is a set of Liouville measure $`0`$.
To prove an upper bound involving the dimension of $`K`$ we combine the methods used to prove Theorems 1 and 2. There exist functions $`\phi _\pm 𝒞^{1,1}(T^{}X)`$ such that, uniformly on compact sets,
$$H_p\phi _\pm \phi _\pm ,\phi _\pm d(\mathrm{\Gamma }_\pm ,)^2,\phi _++\phi _{}d(K,)^2.$$
A local model for the simplest case of one trajectory is given by $`p=\xi _1+x_2\xi _2`$, $`(x,\xi )T^{}^2`$, so that
(2.3)
$$H_p=_{x_1}+x_2_{x_2}\xi _2_{\xi _2},\phi _+=\xi _2^2,\phi _{}=x_2^2,K=\{(t,0;0,0):t\}.$$
A new escape function is given by
(2.4)
$$G\stackrel{\mathrm{def}}{=}\left(\mathrm{log}(Cϵ+\widehat{\phi }_{})\mathrm{log}(Cϵ+\widehat{\phi }_+)\right),ϵMh,M1,$$
where $`\widehat{\phi }_\pm `$ are suitable $`h`$-dependent regularizations of $`\phi _\pm `$.
The logarithmic flattening of the more straightforward escape function $`\phi _{}\phi _+`$ is forced by the requirement that $`G=𝒪(\mathrm{log}(1/h))`$ so that the conjugation used in the proof of Theorem 1 can be applied. However, even for uniformly smooth $`\widehat{\phi }_\pm `$ the regularization of $`G`$ is essentially in the symbolic class $`S_{\frac{1}{2}}`$ and the situation becomes more complicated in general. Nevertheless we obtain the following estimates:
$$_{(x,\xi )}^\alpha H_p^kG=𝒪(ϵ^{\frac{|\alpha |}{2}}),\text{for }|\alpha |+k1\text{, uniformly on compact sets,}$$
and
$$d((x,\xi ),K)^2CϵH_pG1/C.$$
As in the proof of Theorem 1 (but with very different parameters and escape functions) we introduce a conjugated operator,
$$P_{\theta ,t}(h)\stackrel{\mathrm{def}}{=}e^{tG^w}P_\theta (h)e^{tG^w},$$
which now is in an exotic $`\frac{1}{2}`$-class, with the second Planck constant $`\stackrel{~}{h}1/M`$ playing the rôle of the asymptotic parameter. The escape function used here, $`G`$, has compact support.
We now build a second microlocal calculus which combines this exotic class with the one used in the proof of Theorem 2. The first allows us the use of the irregular escape function and the second allows a localization to an $`h`$-neighbourhood of the energy surface. In the new calculus the operator
$$\stackrel{~}{P}_{\theta ,t}=P_{\theta ,t}iMh\stackrel{~}{\mathrm{Op}_h^w}(a),a(x,\xi )\stackrel{\mathrm{def}}{=}\chi \left(\frac{p(x,\xi )}{K_1h}\right)\chi \left(\frac{C_0H_pG(x,\xi )}{K_1h}\right),$$
is globally elliptic (here $`\stackrel{~}{\mathrm{Op}_h^w}`$ describes a second microlocal quantization operator). As in the proof of Theorem 2 the number of eigenvalues of $`P_{\theta ,t}`$, and hence $`P_\theta `$, near $`0`$, is estimated by $`h^n`$ times the volume of the support of $`a`$. A cross-section of that support with a hypersuface transversal to the flow is illustrated in Fig.5. That volume is bounded by $`h^{1+(2n22\nu )/2}`$, where $`\nu >\nu _0`$, $`2\nu _0+1`$ is the dimension of $`K`$. That gives (1.3).
## 3. Preliminaries
### 3.1. Review of semiclassical pseudodifferential calculus
Let $`X`$ be a $`𝒞^{\mathrm{}}`$ manifold which is agrees with $`^n`$ outside a compact set, or more generally
$$X=X_0(^nB(0,R_0))\mathrm{}(^nB(0,R_0)),X_0X.$$
We introduce the usual class of semiclassical symbols on $`X`$:
$$S^{m,k}(T^{}X)=\{a𝒞^{\mathrm{}}(T^{}X\times (0,1]):|_x^\alpha _\xi ^\beta a(x,\xi ;h)|C_{\alpha ,\beta }h^m\xi ^{k|\beta |}\},$$
where outside $`X_0`$ we take the usual $`^n`$ coordinates in this definition. The corresponding class of pseudodifferential operators is denoted by $`\mathrm{\Psi }_h^{m,k}(X)`$, and we have the quantization and symbol maps:
$$\begin{array}{cc}& \mathrm{Op}_h^w:S^{m,k}(T^{}X)\mathrm{\Psi }_h^{m,k}(X)\hfill \\ & \sigma _h:\mathrm{\Psi }_h^{m,k}(X)S^{m,k}(T^{}X)/S^{m1,k1}(T^{}X),\hfill \end{array}$$
with both maps surjective, and the usual properties
(3.3)
$$\begin{array}{c}\sigma _h(AB)=\sigma _h(A)\sigma _h(B),\\ 0\mathrm{\Psi }^{m1,k1}(X)\mathrm{\Psi }^{m,k}(X)\stackrel{\sigma _h}{}S^{m,k}(T^{}X)/S^{m1,k1}(T^{}X)0,\end{array}$$
a short exact sequence, and
$$\sigma _h\mathrm{Op}_h^w:S^{m,k}(T^{}X)S^{m,k}(T^{}X)/S^{m1,k1}(T^{}X),$$
the natural projection map. The class of operators and the quantization map are defined locally using the definition on $`^n`$:
(3.4)
$$\mathrm{Op}_h^w(a)u(x)=\frac{1}{(2\pi h)^n}a(\frac{x+y}{2},\xi )e^{ixy,\xi /h}u(y)𝑑y𝑑\xi ,$$
and we refer to for a detailed discussion. We remark only that when we consider the operators acting on half-densities we can define the symbol map, $`\sigma _h`$, onto
$$S^{m,k}(T^{}X)/S^{m2,k2}(T^{}X),$$
see \[33, Appendix\]. We keep this in mind but for notational simplicity we suppress the half-density notation.
For $`aS^{m,k}(T^{}X)`$ we define
$$\text{ess-supp}_haT^{}XS^{}X,S^{}X\stackrel{\mathrm{def}}{=}(T^{}X0)/_+,$$
where the usual $`_+`$ action is given by multiplication on the fibers: $`(x,\xi )(x,t\xi )`$, as
$$\begin{array}{cc}\hfill \text{ess-supp}_ha=& \mathrm{}\{(x,\xi )T^{}X:ϵ>0,_x^\alpha _\xi ^\beta a(x^{},\xi ^{})=𝒪(h^{\mathrm{}}),d(x,x^{})+|\xi \xi ^{}|<ϵ\}\hfill \\ & \mathrm{}\{(x,\xi )T^{}X0:ϵ>0,_x^\alpha _\xi ^\beta a(x^{},\xi ^{})=𝒪(h^{\mathrm{}}\xi ^{}^{\mathrm{}}),\hfill \\ & d(x,x^{})+1/|\xi ^{}|+|\xi /|\xi |\xi ^{}/|\xi ^{}||<ϵ\}/_+,\hfill \end{array}$$
where the second complement is in $`S^{}X`$. For $`A\mathrm{\Psi }_h^{m,k}(X)`$, then define
$$\mathrm{WF}_h(A)=\text{ess-supp}_ha,A=\mathrm{Op}_h^w(a),$$
noting that, as usual, the definition does not depend on the choice of $`\mathrm{Op}_h^w`$, and
$$\mathrm{Char}(A)=\{WF_h(B):B\mathrm{\Psi }_h^{m,k}(X),\sigma _h(B)=\sigma _h(A)\},$$
where $`\sigma _h`$ is the principal symbol map in (3.3). For
$$u𝒞^{\mathrm{}}((0,1]_h;𝒞^{\mathrm{}}(X)),KX,NP,h_0,u_{C^N(K)}h^P,h<h_0.$$
we define
$$\mathrm{WF}_h(u)=\left(\{\mathrm{Char}(A):A\mathrm{\Psi }^{0,0}(X):Auh^{\mathrm{}}𝒞^{\mathrm{}}((0,1]_h;𝒞^{\mathrm{}}(X))\}\right)^{\mathrm{}},$$
where the complement is taken in $`T^{}XS^{}X`$. Here we will be concerned with a purely semiclassical theory and deal only with compact subsets of $`T^{}X`$.
To illustrate the $`h`$-pseudodifferential calculus at work we prove two simple lemmas which will be used later. We say that $`A\mathrm{\Psi }^{m,k}(X)`$ is elliptic on $`KT^{}X`$ if $`|\sigma (A)_K|>h^m/C`$. This is equivalent to saying
###### Lemma 3.1.
Suppose $`Q\mathrm{\Psi }^{0,m}(X)`$ is elliptic at $`(x_0,\xi _0)`$, $`u_{L^2}=1`$, and $`\mathrm{WF}_h(u)`$ is contained in a sufficiently small neighbourhood of $`(x_0,\xi _0)`$. Then for $`h`$ small enough,
$$Qu_{L^2}1/C.$$
###### Lemma 3.2.
Suppose that $`\psi _j𝒞_\mathrm{b}^{\mathrm{}}(T^{}X)`$, $`\psi _1^2+\psi _2^2=1`$, $`\mathrm{supp}\psi _1\{(x,\xi ):|\xi |C\}`$. Then, there exist $`\mathrm{\Psi }_1\mathrm{\Psi }^{0,\mathrm{}}(X)`$ and $`\mathrm{\Psi }_2\mathrm{\Psi }^{0,0}(X)`$, with principal symbols $`\psi _1`$ and $`\psi _2`$ respectively, such that
$$\mathrm{\Psi }_1^2+\mathrm{\Psi }_1^2=I+R,R\mathrm{\Psi }^\mathrm{},\mathrm{}(X),\mathrm{\Psi }_j^{}=\mathrm{\Psi }_j.$$
###### Proof.
Functional calculus gives
$$(\psi _1^w)^2+(\psi _2^w)^2=I+r_1^w,r_1S^{1,\mathrm{}}(T^{}X),$$
in particular $`r=𝒪(h):H^M(X)H^M(X)`$. If $`h`$ is small enough we put
$$\mathrm{\Psi }_j^1=(1+r_1^w)^{\frac{1}{4}}\psi _j^w(1+r_1^w)^{\frac{1}{4}},$$
so that
$$(\mathrm{\Psi }_1^1)^2+(\mathrm{\Psi }_2^1)^2=I+r_2^w,r_2S^{2,\mathrm{}}(T^{}X),(\mathrm{\Psi }_j^1)^{}=\mathrm{\Psi }_j^1.$$
and we can then proceed by iteration. ∎
The semiclassical Sobolev spaces, $`H_h^s(X)`$ are defined by choosing a globally elliptic, self-adjoint operator, $`A\mathrm{\Psi }^{0,1}(X)`$ (that is an operator satisfying $`\sigma (A)\xi /C`$ everywhere) and putting
$$u_{H_h^s}=A^su_{L^2(X)}.$$
When $`X=^n`$,
$$u_{H_h^s}^2_^nh\xi ^{2s}|u(\xi )|^2𝑑\xi ,u(\xi )\stackrel{\mathrm{def}}{=}_^nu(x)e^{ix,\xi }𝑑x.$$
The following lemma will also be useful:
###### Lemma 3.3.
Suppose that $`P_t`$, $`t(0,\mathrm{})`$, is a family of operators such that
$$P_t:H_h^s(X)H_h^{sm}(X),$$
$$A\mathrm{\Psi }^{0,\mathrm{}}(X),\mathrm{ad}_{P_t}A=𝒪(h):L^2(X)L^2(X),0<h<h_0(t),.$$
Let $`\mathrm{\Psi }_j`$ be as in Lemma 3.2 and suppose that
$$P_t\mathrm{\Psi }_juth\mathrm{\Psi }_ju𝒪(h/t)u,j=1,2,u𝒞_\mathrm{c}^{\mathrm{}}(X).$$
Here the constants in $`𝒪`$ are independent of $`h`$ and $`t`$. Then for $`t>t_01`$ and $`0<h<h_0(t)`$,
$$P_tuthu/2.$$
###### Proof.
We recall from Lemma 3.2 that
(3.5)
$$\mathrm{\Psi }_1v^2+\mathrm{\Psi }_2v^2=v^2+Rv,v=v^2+𝒪(h^{\mathrm{}})v_{H_h^N},$$
and hence with $`v=P_tu`$,
$$\begin{array}{cc}\hfill P_tu^2& =\mathrm{\Psi }_1P_tu^2+\mathrm{\Psi }_2P_tu^2𝒪(h^{\mathrm{}})u^2\hfill \\ & P_t\mathrm{\Psi }_1u^2+P_t\mathrm{\Psi }_2u^2[\mathrm{\Psi }_2,P_t]u^2[\mathrm{\Psi }_2,P_t]u^2\hfill \\ & \mathrm{\hspace{0.33em}2}\left(\mathrm{\Psi }_1P_tu[\mathrm{\Psi }_1,P_t]u^2+\mathrm{\Psi }_2P_tu[\mathrm{\Psi }_2,P_t]u^2\right)^{\frac{1}{2}}𝒪(h^{\mathrm{}})u^2\hfill \\ & P_t\mathrm{\Psi }_1u^2+P_t\mathrm{\Psi }_2u^2\hfill \\ & 2C([\mathrm{\Psi }_1,P_t]u^2+[\mathrm{\Psi }_2,P_t]u^2)P_tu^2/C𝒪(h^{\mathrm{}})u^2\hfill \\ & P_t\mathrm{\Psi }_1u^2+P_t\mathrm{\Psi }_2u^2C^{}h^2u^2P_tu^2/C.\hfill \end{array}$$
We now use the hypothesis of the lemma and (3.5) with $`v=u`$ to obtain
$$\begin{array}{cc}\hfill P_tu^2& t^2h^2(\mathrm{\Psi }_1u^2+\mathrm{\Psi }_2u^2)C^{}h^2u^2P_tu^2/C\hfill \\ & t^2h^2u^2C^{}h^2u^2P_tu^2/C\hfill \end{array}$$
and the lemma follows. ∎
### 3.2. Semiclassical Fourier integral operators.
We now follow and review some aspects of the theory of semiclassical Fourier Integral Operators. We take a point of view which will be used in showing invariance of the second microlocal calculus developed below in §5.
Thus let $`A(t)`$ be a smooth family of pseudodifferential operators,
$$A(t)=\mathrm{Op}_h^w(a(t)),a(t)𝒞^{\mathrm{}}([1,1]_t;S^{0,\mathrm{}}(T^{}X)),$$
such that for all $`t`$, $`\mathrm{WF}_h(A(t))T^{}X`$. We then define a family of operators
(3.8)
$$\begin{array}{c}U(t):L^2(X)L^2(X),\\ hD_tU(t)+U(t)A(t)=0,U(0)=U_0\mathrm{\Psi }_h^{0,0}(X).\end{array}$$
This is an example of a family of $`h`$-Fourier Integral Operators, $`U(t)`$, associated to canonical transformations $`\kappa (t)`$, generated by the Hamilton vector fields $`H_{a_0(t)}`$, where the real valued $`a_0(t)`$ is the $`h`$-principal symbol of $`A(t)`$,
$$\frac{d}{dt}\kappa (t)(x,\xi )=(\kappa (t))_{}(H_{a_0(t)}(x,\xi )),\kappa (0)(x,\xi )=(x,\xi ),(x,\xi )T^{}X.$$
We will often need the Egorov theorem which can be proved directly from this definition: when $`U_0`$ in (3.8) is elliptic (that is $`|\sigma (U_0)|>c>0`$) on $`T^{}X`$, then for $`B\mathrm{\Psi }_h^{m,k}(X)`$
(3.10)
$$\begin{array}{c}\sigma (V(t)BU(t))=(\kappa (t))^{}\sigma (B),\end{array}$$
where $`V(t)`$ is an approximate inverse to $`U(t)`$,
$$V(t)U(t)I,U(t)V(t)I\mathrm{\Psi }_h^\mathrm{},\mathrm{}(T^{}X).$$
The approximate inverse is constructed by taking
$$hD_tV(t)A(t)V(t)=0,V(0)=V_0,V_0U_0I,U_0V_0I\mathrm{\Psi }_h^\mathrm{},\mathrm{}(T^{}X),$$
the existence of $`V_0`$ being guaranteed by the ellipticity of $`U_0`$. The proof of (3.10) follows from writing $`B(t)=V(t)BU(t)`$, so that, in view of the properties of $`V(t)`$,
$$hD_tB(t)[A(t),B(t)]mod\mathrm{\Psi }_h^\mathrm{},\mathrm{},B(0)=B_0.$$
Since the symbol of the commutator is given by $`(h/i)H_{a_0(t)}\sigma (B(t))`$, (3.10) follows directly from the definition of $`\kappa (t)`$.
If $`U=U(1)`$, say, and the graph of $`\kappa (1)`$ is denoted by $`C`$, we conform to the usual notation and write
$$UI_h^0(X\times X;C^{}),C^{}=\{(x,\xi ;y,\eta ):(x,\xi )=\kappa (y,\eta )\},$$
which means that $`U`$ is an $`h`$-Fourier Integral Operator associated to the canonical graph $`C`$. Locally all $`h`$-Fourier Integral Operators associated to canonical graphs are of the form $`U(1)`$ since each local canonical transformation with a fixed point can be deformed to identity, see \[33, Lemma 3.2\] and the proof of Lemma 5.8 below.
Our definitions of pseudo-differential operators and of (the special class of) $`h`$-Fourier Integral Operators were global. It is useful and natural to consider the operators and their properties microlocally. We consider classes of tempered operators:
$$T:𝒞^{\mathrm{}}(X)𝒞^{\mathrm{}}(X),$$
and for any semi-norm $`_1`$ on $`𝒞^{\mathrm{}}(X)`$ there exist a seminorm $`_2`$ on $`𝒞^{\mathrm{}}(X)`$ and a constant $`M_0`$ such that
$$Tu_1=𝒪(h^{M_0})u_2.$$
We remark that since we deal with compact subsets of $`T^{}X`$ here, we could consider operators $`T:𝒟^{}(X)𝒞^{\mathrm{}}(X)`$ in which case we can ask for existence of $`M_0`$ for any two seminorms $`_j`$, $`j=1,2`$.
For open sets, $`VT^{}X`$, $`UT^{}X`$, the operators defined microlocally near $`V\times U`$ are given by equivalence classes of tempered operators given by the relation
$$TT^{}A(TT^{})B=𝒪(h^{\mathrm{}}):𝒟^{}(X)𝒞^{\mathrm{}}(X),$$
for any $`A,B\mathrm{\Psi }_h^{0,0}(X)`$ such that
(3.13)
$$\begin{array}{c}\mathrm{WF}_h(A)\stackrel{~}{V},\mathrm{WF}_h(B)\stackrel{~}{U},\\ \overline{V}\stackrel{~}{V}T^{}X,\overline{U}\stackrel{~}{U}T^{}X,\stackrel{~}{U},\stackrel{~}{V}\text{ open}.\end{array}$$
The equivalence class $`T`$, $`h`$-Fourier Integral Operator associated to a local canonical graph $`C`$ if, again for any $`A`$ and $`B`$ above
$$ATBI^0(X\times X;\stackrel{~}{C}^{}),$$
where $`C`$ needs to be defined only near $`U\times V`$.
We say that $`P=Q`$ microlocally near $`U\times V`$ if $`APBAQB=𝒪_{L^2L^2}(h^{\mathrm{}})`$, where because of the assumed pre-compactness of $`U`$ and $`V`$ the $`L^2`$ norms can be replaced by any other norms. For operator identities this will be the meaning of equality of operators in this paper, with $`U,V`$ specified (or clear from the context). Similarly, we say that $`B=T^1`$ microlocally near $`U\times V`$, if $`BT=I`$ microlocally near $`U\times U`$, and $`TB=I`$ microlocally near $`V\times V`$. More generally, we could say that $`P=Q`$ microlocally on $`WT^{}X\times T^{}X`$ (or, say, $`P`$ is microlocally defined there), if for any $`U,V`$, $`U\times VW`$, $`P=Q`$ microlocally in $`U\times V`$.
If the open sets $`U`$ or $`V`$ in (3.13) are small enough, so that they can be identified with neighbourhoods of points in $`T^{}^n`$, we can use that identification to state that $`T`$ is microlocally defined near, say, $`(m,(0,0))`$, $`mT^{}X`$, $`(0,0)T^{}^n`$.
To give a useful example of this formalism we state a semiclassical version of Egorov’s theorem.
###### Proposition 3.4.
Suppose that $`F`$ is an $`h`$-Fourier integral operator, microlocally defined near $`U\times VT^{}X\times T^{}X`$, and associated to a locally defined canonical transformation $`\kappa :VU`$, and elliptic near $`(\kappa (\rho ),\rho )U\times V`$. Let $`F^1`$ be the microlocal inverse of $`F`$ near $`(\kappa (\rho ),\rho )`$. Then for any $`A\mathrm{\Psi }^{m,k}(X)`$,
(3.14)
$$F^1AF=B\mathrm{\Psi }^{m,k}(X),\sigma _h(B)=\sigma _h(A)\kappa ,$$
microlocally near $`\rho T^{}X`$.
The proof is a localized adaptation of the argument giving (3.10) and the observation recalled above that any canonical transformation with a fixed point (as we can assume that $`\kappa (m)=m`$) can be deformed to the identity.
### 3.3. $`S_{\frac{1}{2}}`$ spaces with two parameters
We define the following symbol class:
(3.15)
$$aS_{\frac{1}{2}}^{m,\stackrel{~}{m},k}(T^{}^n)|_x^\alpha _\xi ^\beta a(x,\xi )|C_{\alpha \beta }h^m\stackrel{~}{h}^{\stackrel{~}{m}}\left(\frac{\stackrel{~}{h}}{h}\right)^{\frac{1}{2}(|\alpha |+|\beta |)}\xi ^{k|\beta |},$$
where in the notation we suppress the dependence of $`a`$ on $`h`$ and $`\stackrel{~}{h}`$. We define the Weyl quantization of $`a`$ in the usual way
$$a^w(x,hD_x)u=\frac{1}{(2\pi h)^n}a(\frac{x+y}{2},\xi )e^{\frac{i}{h}xy,\xi }u(y)𝑑y𝑑\xi ,$$
and the standard results (see ) show that if $`aS_{\frac{1}{2}}^{m,\stackrel{~}{m},k}(T^{}^n)`$ and $`bS_{\frac{1}{2}}^{m^{},\stackrel{~}{m}^{},k^{}}(T^{}^n)`$ then
$$a(x,hD_x)b(x,hD_x)=c(x,hD_x)\text{ with }cS_{\frac{1}{2}}^{m+m^{},\stackrel{~}{m}+\stackrel{~}{m}^{},k+k^{}}(T^{}^n).$$
The presence of the additional parameter $`\stackrel{~}{h}`$ allows us to conclude that
$$c\underset{|\alpha |<M}{}\frac{1}{\alpha !}_\xi ^\alpha aD_x^\alpha bmodS_{\frac{1}{2}}^{m+m^{},\stackrel{~}{m}+\stackrel{~}{m}^{}M,k+k^{}M}(T^{}^n),$$
that is, we have a symbolic expansion in powers of $`\stackrel{~}{h}`$. We could also consider an expansion in the Weyl quantization – see (3.21).
We denote our class of operators by $`\mathrm{\Psi }_{\frac{1}{2}}^{m,\stackrel{~}{m},k}(T^{}^n)`$. For simplicity we will only state the characterization à la Beals for a simpler class of symbols:
###### Lemma 3.5.
Suppose that $`A:𝒮(^n)𝒮^{}(^n)`$. Then $`A=\mathrm{Op}_h^w(a)`$ with
(3.16)
$$_x^\alpha _\xi ^\beta a=𝒪(h^m\stackrel{~}{h}^{\stackrel{~}{m}})\left(\frac{\stackrel{~}{h}}{h}\right)^{\frac{1}{2}(|\alpha |+|\beta |)},$$
if and only if for any sequence $`\{\mathrm{}_j\}_{j=1}^N`$ of linear functions on $`T^{}^n`$ we have
$$\mathrm{ad}_{\mathrm{Op}_h^w(\mathrm{}_1)}\mathrm{}\mathrm{ad}_{\mathrm{Op}_h^w(\mathrm{}_N)}Au_{L^2(^n)}Ch^{m+N/2}\stackrel{~}{h}^{\stackrel{~}{m}+N/2}u_{L^2(^n)},$$
for any $`u𝒮(^n)`$.
###### Proof.
We can assume that $`m=\stackrel{~}{m}=0`$. The statement follows from the proof in \[10, Chapter 8\] and a rescaling:
$$(\stackrel{~}{x},\stackrel{~}{\xi })=(\stackrel{~}{h}/h)^{\frac{1}{2}}(x,\xi ).$$
In fact, we define the following unitary operator on $`L^2(^n)`$:
$$U_{h,\stackrel{~}{h}}u(\stackrel{~}{x})=(\stackrel{~}{h}/h)^{\frac{n}{4}}u((h/\stackrel{~}{h})^{\frac{1}{2}}\stackrel{~}{x}),$$
for which we can check that
$$a(x,hD_x)=U_{h,\stackrel{~}{h}}^1a_{h,\stackrel{~}{h}}(\stackrel{~}{x},\stackrel{~}{h}D_{\stackrel{~}{x}})U_{h,\stackrel{~}{h}},a_{h,\stackrel{~}{h}}(\stackrel{~}{x},\stackrel{~}{\xi })=a((h/\stackrel{~}{h})^{\frac{1}{2}}(\stackrel{~}{x},\stackrel{~}{\xi })).$$
Clearly $`a`$ satisfies (3.16) if and only if $`a_{h,\stackrel{~}{h}}𝒞_\mathrm{b}^{\mathrm{}}(T^{}^n)`$. The Beals condition for $`\stackrel{~}{h}`$-pseudodifferential operators is
$$\mathrm{ad}_{\stackrel{~}{\mathrm{}}_1(\stackrel{~}{x},\stackrel{~}{h}D_{\stackrel{~}{x}})}\mathrm{}\mathrm{ad}_{\stackrel{~}{\mathrm{}}_N(\stackrel{~}{x},\stackrel{~}{h}D_{\stackrel{~}{x}})}a_{h,\stackrel{~}{h}}(\stackrel{~}{x},\stackrel{~}{h}D_{\stackrel{~}{x}})u_{L^2}C\stackrel{~}{h}^Nu_{L^2}.$$
But this is the condition in the lemma since we should take
$$\stackrel{~}{\mathrm{}}_j=(\mathrm{}_j)_{h,\stackrel{~}{h}}=(\stackrel{~}{h}/h)^{\frac{1}{2}}\mathrm{}_j,$$
and this completes the proof. ∎
We remark that the proof given in ,\[10, Chapter 8\] is recalled, in a more complicated setting, in the proof of Proposition 5.2 below. It is based on showing that
$$\alpha ^{2n},(^\alpha a)(x,D)u_{L^2}C_\alpha u_{L^2}\beta ^{2n},sup|^\beta a|C_\beta ^{},$$
which in the setting presented in Lemma 3.5 becomes
(3.20)
$$\begin{array}{c}\alpha ^{2n},(^\alpha a)(x,hD)u_{L^2}C_\alpha (\stackrel{~}{h}/h)^{|\alpha |/2}u_{L^2}\\ \\ \beta ^{2n},sup|^\beta a|C_\beta ^{}(\stackrel{~}{h}/h)^{|\alpha |/2}.\end{array}$$
We will also need the following application of the semi-classical calculus:
###### Lemma 3.6.
Suppose that $`^\alpha a`$, $`^\alpha b=𝒪_\alpha ((\stackrel{~}{h}/h)^{|\alpha |/2}),`$ and that $`c^w(x,hD)=a^w(x,hD)b^w(x,hD)`$. Then
(3.21)
$$c(x,\xi )=\underset{k=0}{\overset{N}{}}\frac{1}{k!}\left(\frac{ih}{2}\sigma (D_x,D_\xi ;D_y,D_\eta )\right)^ka(x,\xi )b(y,\eta )_{x=y,\xi =\eta }+e_N(x,\xi ),$$
where for some $`M`$
(3.22)
$$\begin{array}{cc}& |^\alpha e_N|C_Nh^{N+1}\times \hfill \\ & \underset{\alpha _1+\alpha _2=\alpha }{}\underset{\genfrac{}{}{0pt}{}{(x,\xi )T^{}^n}{(y,\eta )T^{}^n}}{sup}\underset{|\beta |M,\beta ^{2n}}{sup}\left|(h^{\frac{1}{2}}_{(x,\xi ;y,\eta )})^\beta (i\sigma (D)/2)^{N+1}^{\alpha _1}a(x,\xi )^{\alpha _2}b(y,\eta )\right|,\hfill \end{array}$$
where $`\sigma (D)=\sigma (D_x,D_\xi ;D_y,D_\eta ).`$
###### Proof.
This follows from from the standard estimates of symbolic calculus (see \[10, Proposition 7.6\]): suppose that $`A(D)`$ is a non-degenerate real quadratic form. Then there exists $`M`$ such that
$$|^\alpha \mathrm{exp}(iA(D))a(x,\xi )|C\underset{|\beta |M}{}\underset{(x,\xi )T^{}^n}{sup}|^{\alpha +\beta }a(x,\xi )|.$$
We observe that a rescaling $`\stackrel{~}{x}=x/\sqrt{s}`$, $`s>0`$, shows that
$$|^\alpha \mathrm{exp}(isA(D))a(x,\xi )|C\underset{|\beta |M}{}\underset{(x,\xi )T^{}^n}{sup}|^\alpha (\sqrt{s})^\beta a(x,\xi )|.$$
To obtain an expansion we use the Taylor expansion:
$$\mathrm{exp}(ihA(D))=\underset{k=0}{\overset{N}{}}\frac{(ihA(D))^k}{k!}+\frac{1}{N!}_0^1(1t)^N\mathrm{exp}(ithA(D))(ihA(D))^{N+1}𝑑t.$$
In the notation of the lemma and with $`A(D)=\sigma (D_x,D_\xi ;D_y,D_\eta )/2`$,
$$c(x,\xi )=\mathrm{exp}(iA(D))a(x,\xi )b(y,\eta )_{x=y,\eta =\xi },$$
and the lemma follows. ∎
As a particular consequence we notice that if $`aS_{\frac{1}{2}}^{0,0,\mathrm{}}(T^{}^n)`$ and $`bS^{0,\mathrm{}}(T^{}^n)`$ then
$$a^w(x,hD)b^w(x,hD)=c^w(x,hD),c(x,\xi )=$$
$$\underset{k=0}{\overset{N}{}}\frac{1}{k!}\left(ih\sigma (D_x,D_\xi ;D_y,D_\eta )\right)^ka(x,\xi )b(y,\eta )_{x=y,\xi =\eta }+𝒪(h^{\frac{N+1}{2}}\stackrel{~}{h}^{\frac{N+1}{2}}),$$
and the usual pseudodifferential calculus allows a remainder improvement to
$$𝒪(h^{\frac{N+1}{2}}\stackrel{~}{h}^{\frac{N+1}{2}}\xi ^{\mathrm{}}).$$
### 3.4. One parameter groups of elliptic operators.
We recall a special case of a result of Bony and Chemin \[4, Théoreme 6.4\]. Let $`m(x,\xi )`$ be an order function in the sense of :
(3.23)
$$m(x,\xi )Cm(y,\eta )(xy,\xi \eta )^N.$$
The class of symbols, $`S(m)`$, corresponding to $`m`$ is defined as
$$aS(m)|_x^\alpha _\xi ^\beta a(x,\xi )|C_{\alpha \beta }m(x,\xi ).$$
If $`m_1`$ and $`m_2`$ are order functions in the sense of (3.23), and $`a_jS(m_j)`$ then (we put $`h=1`$ here),
$$a_1^w(x,D)a_2^w(x,D)=b^w(x,D),bS(m_1m_2),$$
with $`b`$ given by the usual formula,
(3.24)
$$\begin{array}{cc}\hfill b(x,\xi )& =a_1\mathrm{\#}a_2(x,\xi )\hfill \\ & \stackrel{\mathrm{def}}{=}\mathrm{exp}(i\sigma (D_{x^1},D_{\xi ^1};D_{x^2},D_{\xi ^2})/2)a_1(x^1,\xi ^1)a_2(x^2,\xi ^2)_{x^1=x^2=x,\xi ^1=\xi ^2=\xi }.\hfill \end{array}$$
A special case of \[4, Théoreme 6.4\] gives
###### Proposition 3.7.
Let $`m`$ be an order function in the sense of (3.23) and suppose that $`G𝒞_\mathrm{c}^{\mathrm{}}(T^{}^n;)`$ satisfies
(3.25)
$$G(x,\xi )\mathrm{log}m(x,\xi )=𝒪(1),_x^\alpha _\xi ^\beta G(x,\xi )=𝒪(1),|\alpha |+|\beta |1.$$
Then
(3.26)
$$\mathrm{exp}(tG^w(x,D))=B_t^w(x,D),B_tS(m^t).$$
Here $`\mathrm{exp}(tG^w(x,D))`$ is constructed using spectral theory of bounded self-adjoint operators. The estimates on $`B_tS(m^t)`$ depend only on the constants in (3.25) and in (3.23). In particular they are independent of the support of $`G`$.
In Appendix at the end of the paper we give a simple direct proof of this proposition. We should stress that the main difficulties in came from considering general Weyl calculi of pseudodifferential opearators. Here we need only the case of the simplest metric $`g=dx^2+d\xi ^2`$.
### 3.5. Review of complex scaling
We very briefly recall the procedure described in . It follows the long tradition of the complex scaling method – see for the presentation for compactly supported perturbations and references to earlier work.
Let $`\mathrm{\Gamma }_\theta ^n`$ be a totally real contour with the following properties:
(3.30)
$$\begin{array}{c}\mathrm{\Gamma }_\theta B_^n(0,R_0)=B_^n(0,R_0),\\ \mathrm{\Gamma }_\theta ^nB_^n(0,2R_0)=e^{i\theta }^n^nB_^n(0,2R_0),\\ \mathrm{\Gamma }_\theta =\{x+if_\theta (x):x^n\},_x^\alpha f_\theta (x)=𝒪_\alpha (\theta ).\end{array}$$
The contour can be considered as a deformation of the manifold $`X`$ as nothing is being done in the compact region. The operator $`P`$ defines a dilated operator:
$$P_\theta \stackrel{\mathrm{def}}{=}P_{\mathrm{\Gamma }_\theta },P_\theta u=\stackrel{~}{P}(\stackrel{~}{u})_{\mathrm{\Gamma }_\theta },$$
where $`\stackrel{~}{P}`$ is the holomorphic continuation of the operator $`P`$, and $`\stackrel{~}{u}`$ is an almost analytic extension of $`u𝒞_\mathrm{c}^{\mathrm{}}(\mathrm{\Gamma }_\theta )`$ (here we are only concerned with $`\mathrm{\Gamma }_\theta B_^n(0,R_0)`$).
For $`\theta `$ fixed, the scaled operator, $`P_\theta `$, is uniformly elliptic in $`\mathrm{\Psi }^{0,2}(X)`$ outside a compact set (see (4.11) below) and hence the resolvent, $`(P_\theta z)^1`$, is meromorphic for $`zD(0,1/C)`$. We can also take $`\theta `$ to be $`h`$ dependent and the same statement holds for $`zD(0,\theta /C)`$. The spectrum of $`P_\theta `$ in $`zD(0,\theta /C)`$ is independent of $`\theta `$ and consists of quantum resonances of $`P`$ which are defined as the poles of the meromorphic continuation of
$$(Pz)^1:𝒞_\mathrm{c}^{\mathrm{}}(X)𝒞^{\mathrm{}}(X).$$
In fact, that is one of the ways of defining resonances, and in this paper we will be estimating the number of eigenvalues of $`P_\theta `$.
## 4. Resonance free regions under the non-trapping assumption
### 4.1. Estimates using weight functions
We follow the presentation given in \[9, §4.2\] and inspired by many previous works, including .
Let us suppose that $`P_\theta \mathrm{\Psi }^{0,2}(X)`$ (we identify $`X`$ and $`X_\theta `$ here) is a complex scaled operator with $`\theta =M_1h\mathrm{log}(1/h)`$. We choose $`ϵ`$
(4.1)
$$ϵM_2h\mathrm{log}\frac{1}{h},$$
where $`M_2>M_1`$ is a large constant to be fixed later.
Let $`G𝒞_\mathrm{c}^{\mathrm{}}(T^{}X)`$ and define
$$P_{ϵ,\theta }\stackrel{\mathrm{def}}{=}e^{ϵG/h}P_\theta e^{ϵG/h}=e^{\frac{ϵ}{h}\mathrm{ad}_G}P_\theta \underset{0}{\overset{\mathrm{}}{}}\frac{ϵ^k}{k!}(\frac{1}{h}\mathrm{ad}_G)^k(P_\theta ),G=G^w(x,hD).$$
We note that the assumption on $`ϵ`$ and the boundedness of $`\mathrm{ad}_G/h`$ show that the expansion makes sense. The operators $`\mathrm{exp}(ϵG/h)`$ are pseudo-differential in an exotic class $`S_\delta ^{C_2}`$ for any $`\delta >0`$ (see ) but that is not relevant here.
Using the same letters for operators and and the corresponding symbols, we see that
$$P_{ϵ,\theta }=P_\theta iϵ\{P_\theta ,G\}+𝒪(ϵ^2)=p_\theta iϵ\{p_\theta ,G\}+𝒪(h+ϵ^2),$$
so that
$$\begin{array}{cc}& ReP_{ϵ,\theta }=Rep_\theta +ϵ\{Imp_\theta ,G\}+𝒪(h+ϵ^2)=Rep_\theta +𝒪(h+\theta ϵ+ϵ^2),\hfill \\ & ImP_{ϵ,\theta }=Imp_\theta ϵ\{Rep_\theta ,G\}+𝒪(h+ϵ^2).\hfill \end{array}$$
We now make the following assumption: for a fixed $`\delta >0`$
(4.2)
$$|Rep_\theta |<\delta Imp_\theta +ϵH_pGc_0ϵ.$$
Now let $`\psi _1,\psi _2𝒞_\mathrm{b}^{\mathrm{}}(T^{}X)`$ be two functions satisfying
$$\psi _1^2+\psi _2^2=1,\psi _1_{|Rep_\theta |<\delta /2}1,\mathrm{supp}\psi _1\{|Rep_\theta |<\delta \}.$$
Lemma 3.2 gives two selfadjoint operators $`\mathrm{\Psi }_1`$ and $`\mathrm{\Psi }_2`$ with principal symbols $`\psi _1`$ and $`\psi _2`$ respectively, such that
$$\mathrm{\Psi }_1^2+\mathrm{\Psi }_1^2=I+R,R=𝒪(h^{\mathrm{}}):H^M(X)H^M.$$
We then write $`P_{ϵ,\theta }=A_{ϵ,\theta }+iB_{ϵ,\theta }`$, where
$$A_{ϵ,\theta }=\frac{1}{2}(P_{ϵ,\theta }+P_{ϵ,\theta }^{}),B_{ϵ,\theta }=\frac{1}{2i}(P_{ϵ,\theta }P_{ϵ,\theta }^{}).$$
The principal symbol of $`B_{ϵ,\theta }`$ is given by $`Imp_\theta ϵH_pG`$ and on the essential support of $`\mathrm{\Psi }_1`$ it is bounded below by $`c_0ϵh`$. Hence the sharp Gårding inequality (see \[10, Theorem 7.12\]) implies that for $`h`$ small enough
$$\begin{array}{cc}\hfill P_{ϵ,\theta }\mathrm{\Psi }_1u\mathrm{\Psi }_1u& |P_{ϵ,\theta }\mathrm{\Psi }_1u,\mathrm{\Psi }_1u||ImP_{ϵ,\theta }\mathrm{\Psi }_1u,\mathrm{\Psi }_1u|\hfill \\ & =B_{ϵ,\theta }\mathrm{\Psi }_1u,\mathrm{\Psi }_1u\frac{ϵ}{C}\mathrm{\Psi }_1u^2,\hfill \end{array}$$
and hence
$$P_{ϵ,\theta }\mathrm{\Psi }_1u\frac{ϵ}{C}\mathrm{\Psi }_1u.$$
On the support of $`\psi _2`$ the operator $`A_{ϵ,\theta }`$ is elliptic and by Lemma 3.1,
$$P_{ϵ,\theta }\mathrm{\Psi }_2u\frac{1}{C}\mathrm{\Psi }_2u𝒪(h^{\mathrm{}})u.$$
We conclude from Lemma 3.3 that
$$P_{ϵ,\theta }u\frac{ϵ}{C}u.$$
This shows that the conjugated operator has no spectrum in $`D(0,ϵ/(2C))`$.
### 4.2. Construction of an escape function
Using the results of §4.1 all we need to do is to construct $`G`$ so that (4.2) holds. For that we modify a standard argument with the presentation borrowed in part from \[36, Sect.4\].
We recall that (1.14) implies the same condition with $`p^1([ϵ_0,ϵ_0])`$ for some $`ϵ_0>0`$. That follows from the compactness of the trapped set in $`p^1([\delta ,\delta ])`$ – see \[11, Appendix\] for a detailed discussion.
Let us now fix $`R`$ a large parameter. We will define $`G_\rho 𝒞_\mathrm{c}^{\mathrm{}}(T^{}X)`$, a local escape function supported in a neighbourhood of the bicharacteristic segment
$$I_\rho =\{\mathrm{exp}(tH_p)(\rho ):t[T,T]\},$$
and which satisfies $`H_pG_\rho 1`$ on the part of $`I_\rho `$ lying over
(4.3)
$$K^{}=\{\rho ^{}T^{}X:|x(\rho ^{})|R\}$$
For that, let $`\mathrm{\Gamma }`$ be a hypersurface through $`\rho `$ which is transversal to $`H_p`$. Then there is a neighbourhood $`U_\rho `$ of $`\rho `$, such that
$$V_\rho =\{\mathrm{exp}(t(U_\rho \mathrm{\Gamma })):t(T1,T+1)\}p^1([ϵ_0/2,ϵ_0/2]),$$
is a neighbourhood of $`I_\rho `$. That neighbourhood can be identified with a product,
$$V_\rho (T1,T+1)\times (U_\rho \mathrm{\Gamma }),$$
and, in this identification, we will choose $`T`$ and $`0<\alpha <1`$ so that
$$(((T1,\alpha T)(\alpha T,T+1))\times (U_\rho \mathrm{\Gamma })))K^{}=\mathrm{}.$$
We now need the following elementary
###### Lemma 4.1.
For any $`0<\alpha <1/2`$ and $`T>0`$ there exist as function $`\chi =\chi _{T,\alpha }𝒞^{\mathrm{}}(;)`$ such that
$$\chi (t)=\{\begin{array}{cc}0& |t|>T\\ t& |t|<\alpha T\end{array},\chi ^{}(t)2\alpha .$$
###### Proof.
The piecewise linear function
$$\chi _\mathrm{\#}(t)=\{\begin{array}{ccc}0& & |t|>T\\ t& & |t|<\alpha T\\ \pm \alpha (Tt)/(1\alpha )& & \alpha T\pm tT\end{array}$$
satisfies $`\chi _{\mathrm{}}^{}{}_{}{}^{}\alpha /(1\alpha )>2\alpha `$ wherever the derivative is defined. A regularization of this function gives $`\chi _{T,\alpha }`$. ∎
Now let $`\varphi _\rho 𝒞_\mathrm{c}^{\mathrm{}}(U_\rho \mathrm{\Gamma })`$ be identically $`1`$ near $`\rho `$, and let $`\chi _T`$ be given by the lemma. Using the product coordinates, we can think of $`\varphi _\rho `$, $`t`$, and hence $`\chi (t)`$, as functions on $`T^{}X`$. The functions $`\varphi _\rho `$ and $`\chi _T(t)`$ have compact support in $`V_\rho `$. Let
$$\psi 𝒞_\mathrm{c}^{\mathrm{}}((ϵ_0,ϵ_0)),\psi _{[ϵ_0/2,ϵ_0/2]}1,$$
and put
(4.4)
$$G_\rho =\chi _T(t)\varphi _\rho \psi (p),G_\rho 𝒞_\mathrm{c}^{\mathrm{}}(V_\rho ).$$
so that
(4.5)
$$H_pG_\rho =\chi _T^{}\varphi _\rho \psi (p),$$
satisfies
$`H_pG_\rho =1`$ on $`V_\rho \{|x|<R\}`$ and $`H_pG_\rho 2\alpha `$ everywhere.
Now let $`KT^{}X`$ be the compact set
(4.6)
$$K=\{\rho p^1([ϵ_0/3,ϵ_0/3]):|x(\rho )|R/2\}.$$
Since $`K`$ is compact, applying the previous argument for every $`\rho K`$ gives a $`U_\rho `$, and a $`U_\rho ^{}U_\rho `$ on which $`\varphi _\rho =1`$. Since $`\{U_\rho ^{}:\rho K\}`$ covers $`K`$, the compactness of $`K`$ shows that we can pass to a finite subcover, $`\{U_{\rho _j}^{}:j=1,\mathrm{},N\}`$. We let
(4.7)
$$G=\underset{j=1}{\overset{N}{}}G_{\rho _j}.$$
The construction of $`G_{\rho _j}`$’s now shows that by choosing $`\alpha `$ small enough (depending on the maximal number of support overlaps we obtain
(4.8)
$$H_pG(\rho )1,\rho p^1((ϵ_0/2,ϵ_0/2))\{|x(\rho )|<R\}\text{and}H_pG(\rho )\delta ,\rho T^{}X.$$
### 4.3. Resonance free region
We now want to choose the scaling so that (4.2) holds with $`G`$ satisfying (4.8). Once that is done the results of §4.1 will give Theorem 3.
For that we choose the complex scaling so that
(4.11)
$$\begin{array}{c}Imp_\theta (x,\xi )\theta \text{when }|p(x,\xi )|ϵ_0\text{ and }|x|R\text{,}\\ Imp_\theta <C_1\theta \text{when }|p(x,\xi )|ϵ_0\text{,}\end{array}$$
where $`R`$ is independent of $`\theta `$. With $`ϵ=M_2h\mathrm{log}(1/h)`$ we now choose $`\theta =M_1h\mathrm{log}(1/h)`$ such that
$$M_1<M_2/C_1,\delta M_2<M_1,$$
where $`C_1`$ comes from (4.11) and $`\delta `$ comes from (4.8). Since we can choose $`\delta `$ as small as we want this can certainly be arranged leading to (4.2).
For completeness we include a quantitative corollary of Theorem 1 from :
Theorem 1.Suppose that the assumptions of Theorem 1 are satisfied and that $`(P_\theta z)^1`$ is the scaled resolvent defined for $`0<\theta <2Mh\mathrm{log}h`$, $`M1`$. Then for $`0<h<h_0(M)`$ we have
$$(P_\theta z)^1_{L^2(\mathrm{\Gamma }_\theta )L^2(\mathrm{\Gamma }_\theta )}=C\frac{\mathrm{exp}(C|Imz|/h)}{h},zD(0,Mh\mathrm{log}(1/h)).$$
Theorem 1 is essentially optimal as shown by the well known one dimensional result going back to Regge (see for a proof and references): if $`V𝒞^N([a,b])`$ is extended by $`0`$ to a potential on $``$, and
$$V(x)\{\begin{array}{cc}(xa)^p\hfill & xa+\hfill \\ (xb)^q\hfill & xb\hfill \end{array},p,q<N,$$
then the scattering poles for $`\mathrm{\Delta }+V(x)`$ are given at high energies by the sequence
$$\lambda _k=\frac{\pi k}{ba}i\alpha \mathrm{log}|k|+𝒪(1),k,\alpha =\frac{p+q+4}{2(ba)}.$$
The semiclassical resonances, $`z_k(h)`$, of $`h^2\mathrm{\Delta }+h^2V(x)`$ are related to these scattering poles by the formula $`z_k(h)=h^2\lambda _k^2`$. Hence
$$Rez_k(h)1Imz_k(h)h\mathrm{log}(1/h).$$
## 5. Second microlocal calculus associated to a hypersurface
To obtain Theorem 2 we need to localize to an $`h`$-size neighbourhood of the energy surface
$$\mathrm{\Sigma }\stackrel{\mathrm{def}}{=}\{(x,\xi )T^{}X:p(x,\xi )=0\}.$$
That means that we have to work with functions of the form
$$a(x,\xi ;h)=\psi (p(x,\xi )/h).$$
The usual quantization procedure (the passage from symbols to pseudodifferential operators) is prohibited as
$$_{x,\xi }^\alpha ah^{|\alpha |}.$$
The troublesome symbols have a special form and we can construct a calculus which includes them by straightening $`\mathrm{\Sigma }`$ locally by means of canonical transformations. That means moving $`\mathrm{\Sigma }`$ to
$$\mathrm{\Sigma }_0=\{\xi _1=0\}.$$
We may then localize to rectangles in the $`(x_1,\xi _1)`$-space of length $`1`$ in $`x_1`$ and of length $`h`$ in $`\xi _1`$. This amounts to a form of semiclassical second microlocalization. The presentation here is essentially self-contained and we refer to for pointers to the literature.
For Theorem 2 we need even more singular calculus related to the $`\mathrm{\Psi }_{\frac{1}{2}}`$ calculus described in §3.3.
We assume that $`\mathrm{\Sigma }`$ is a compact $`𝒞^{\mathrm{}}`$ hypersurface in $`T^{}X`$. Since the delicate constructions will only be used in a compact set and since we are working in the $`𝒞^{\mathrm{}}`$ category this creates no restrictions.
### 5.1. Basic properties
To construct the calculus, let $`\mathrm{\Sigma }T^{}X`$ be a $`𝒞^{\mathrm{}}`$ compact hypersurface. We consider a class of symbols associated to $`\mathrm{\Sigma }`$, a multiindex
$$𝔭\stackrel{\mathrm{def}}{=}(m,\stackrel{~}{m},k_1,k_2),$$
and depending on two small parameters,
$$0<h<\stackrel{~}{h}.$$
(5.1) aSΣ,δ𝔭(TX){near Σ V1Vl1W1Wl2Hqka=𝒪(hmδl1l2h~m~+δl1+l2(h~/h)qk2+δk(1δ)l2),whereV1,,Vl1 are vector fields
tangent to Σ and W1,Wl2 are
any vectorfields away from Σ xαξβa(x,ξ;h)=𝒪(hmk2δ(|α|+|β|)h~m~+k2+δ(|α|+|β|)ξk1|β|).𝑎superscriptsubscript𝑆Σ𝛿𝔭superscript𝑇𝑋casesnear Σ subscript𝑉1subscript𝑉subscript𝑙1subscript𝑊1subscript𝑊subscript𝑙2superscriptsubscript𝐻𝑞𝑘𝑎absent𝒪superscript𝑚𝛿subscript𝑙1subscript𝑙2superscript~~𝑚𝛿subscript𝑙1subscript𝑙2superscriptdelimited-⟨⟩~𝑞subscript𝑘2𝛿𝑘1𝛿subscript𝑙2wheresubscript𝑉1subscript𝑉subscript𝑙1 are vector fields
tangent to Σ and subscript𝑊1subscript𝑊subscript𝑙2 are
any vectorfields missing-subexpressionaway from Σ superscriptsubscript𝑥𝛼superscriptsubscript𝜉𝛽𝑎𝑥𝜉absent𝒪superscript𝑚subscript𝑘2𝛿𝛼𝛽superscript~~𝑚subscript𝑘2𝛿𝛼𝛽superscriptdelimited-⟨⟩𝜉subscript𝑘1𝛽a\in S_{\Sigma,\delta}^{{\mathfrak{p}}}(T^{*}X)\Longleftrightarrow\left\{\begin{array}[]{l}\text{\text@underline{near $\Sigma$} }\ \ \ V_{1}\cdots V_{l_{1}}W_{1}\cdots W_{l_{2}}H_{q}^{k}a=\\
{\mathcal{O}}\left(h^{-m-\delta l_{1}-l_{2}}\tilde{h}^{-{\widetilde{m}}+\delta l_{1}+l_{2}}\langle(\tilde{h}/h)q\rangle^{k_{2}+\delta k-(1-\delta)l_{2}}\right)\,,\\
\text{where}\ V_{1},\cdots,V_{l_{1}}\ \text{ are vector fields
tangent to $\Sigma$ and }\\
W_{1},\cdots W_{l_{2}}\ \text{ are
any vectorfields }\\
\\
\text{\text@underline{away from $\Sigma$} }\ \ \partial_{x}^{\alpha}\partial_{\xi}^{\beta}a(x,\xi;h)=\\
{\mathcal{O}}(h^{-m-k_{2}-\delta(|\alpha|+|\beta|)}\tilde{h}^{-{\widetilde{m}}+k_{2}+\delta(|\alpha|+|\beta|)}\langle\xi\rangle^{k_{1}-|\beta|})\,.\end{array}\right.
Here $`q`$ is any defining function of $`\mathrm{\Sigma }`$, that is a function which vanishes simply on $`\mathrm{\Sigma }`$. To assure invariance we need
###### Lemma 5.1.
Definition (5.1) is independent of the choice of $`q`$ and the vector fields applied to $`a`$ can be taken in any order.
###### Proof.
The independence of the order of vector fields follows from the fact that $`[V_\mathrm{}_j,H_q]`$ is a vector field tangent to $`\mathrm{\Sigma }`$. To see independence of $`q`$ we note that $`H_{uq}=uH_q+qH_u`$. The vector field $`H_u`$ can be considered as an arbitrary vector field and $`qH_u`$ is tangent to $`\mathrm{\Sigma }`$. Hence the application of the second term is estimated by
$$|qH_ua|C\mathrm{min}(|q|h^1qh^1^{(1\delta )},h^\delta )C^{}h^1q^\delta ,$$
and this estimate can be iterated. Hence we have the same estimates for $`q`$ replaced by $`uq`$, $`u0`$ near $`\mathrm{\Sigma }`$. ∎
The symbol classes in (5.1) are best understood in the simple case when $`q=\xi _1`$, say. The condition for $`|\xi |C`$, means that, with $`\xi =(\xi _1,\xi ^{})`$,
(5.2)
$$_{x_1}^k_x^{}^\alpha _\xi ^{}^\beta (\xi _1_{\xi _1})^\mathrm{}_1_{\xi _1}^\mathrm{}_2a=𝒪\left(h^m\stackrel{~}{h}^{\stackrel{~}{m}}(\stackrel{~}{h}/h)^{\delta (|\alpha |+|\beta |)+\delta \mathrm{}_1+\mathrm{}_2}(\stackrel{~}{h}/h)\xi _1^{k_2+\delta k(1\delta )\mathrm{}_2)}\right),$$
or, if we eliminate the vector fields vanishing at $`\xi _1=0`$,
(5.3)
$$_{x_1}^k_x^{}^\alpha _\xi ^{}^\beta ((h/\stackrel{~}{h})_{\xi _1})^pa=𝒪\left(h^m\stackrel{~}{h}^{\stackrel{~}{m}}(\stackrel{~}{h}/h)^{\delta (|\alpha |+|\beta |)}(\stackrel{~}{h}/h)\xi _1^{k_2(1\delta )p+\delta k}\right).$$
We used the fact that if $`|\xi _1|C`$ then
$$\left((\stackrel{~}{h}/h)\xi _1\right)^{\delta 1}\left((\stackrel{~}{h}/h)\xi _1\right)^1(\stackrel{~}{h}/h)^\delta ,$$
to eliminate the need for the $`\xi _1_{\xi _1}`$. The advantage of the formulation in (5.1) (and (5.2)) is the geometric invariance. In explicit coordinates (5.3) is however more transparent.
Since it is sufficient for our purposes, and for simplicity of presentation we will consider the case of $`\delta =1/2`$ only.
The second microlocalization associates to the space of symbols $`S_{\mathrm{\Sigma },\frac{1}{2}}^𝔭`$ a space of operators, $`\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^𝔭`$, defined in §5.4 below. The basic properties of these spaces are described in the following
###### Theorem 4.
Let us define two multiindices,
$$𝔭=(m,\stackrel{~}{m},k_1,k_2),𝔭^{}=(m,\stackrel{~}{m}1,k_11,k_2).$$
With the definitions of $`S_{\mathrm{\Sigma },\delta }^𝔭`$ above and $`\mathrm{\Psi }_{\mathrm{\Sigma },\delta }^𝔭`$ in §5.4 below, there exist maps
$`\mathrm{Op}_{\mathrm{\Sigma },h,\stackrel{~}{h}}:S_{\mathrm{\Sigma },\frac{1}{2}}^𝔭(T^{}X)\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^𝔭(X)`$
$`\sigma _{\mathrm{\Sigma },h,\stackrel{~}{h}}:\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^𝔭(X)S_{\mathrm{\Sigma },\frac{1}{2}}^𝔭(T^{}X)/S_{\mathrm{\Sigma },\frac{1}{2}}^𝔭^{}(T^{}X)`$
such that
(5.4)
$$\sigma _{\mathrm{\Sigma },h,\stackrel{~}{h}}(AB)=\sigma _{\mathrm{\Sigma },h,\stackrel{~}{h}}(A)\sigma _{\mathrm{\Sigma },h,\stackrel{~}{h}}(B)$$
$$0\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^𝔭^{}(X)\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^𝔭(X)\stackrel{\sigma _{\mathrm{\Sigma },h,\stackrel{~}{h}}}{}S_{\mathrm{\Sigma },\frac{1}{2}}^𝔭(T^{}X)/S_{\mathrm{\Sigma },\frac{1}{2}}^𝔭^{}(T^{}X)0$$
is a short exact sequence and
$$\sigma _{\mathrm{\Sigma },h,\stackrel{~}{h}}\mathrm{Op}_{\mathrm{\Sigma },h,\stackrel{~}{h}}:S_{\mathrm{\Sigma },\frac{1}{2}}^𝔭(T^{}X)S_{\mathrm{\Sigma },\frac{1}{2}}^𝔭(T^{}X)/S_{\mathrm{\Sigma },\frac{1}{2}}^𝔭^{}(T^{}X)$$
is the natural projection map. If
$$aS_{\mathrm{\Sigma },\frac{1}{2}}^𝔭(T^{}X),d(\mathrm{supp}a,\mathrm{\Sigma })1/C$$
then
$$aS_{\frac{1}{2}}^{m+k_2,\stackrel{~}{m}k_2,k_1}(T^{}X),\mathrm{Op}_{\mathrm{\Sigma },h,\stackrel{~}{h}}(a)=\mathrm{Op}_h^w(a)\mathrm{\Psi }_{\frac{1}{2}}^{m+k_2,\stackrel{~}{m}k_2,k_1}(X),$$
where $`\mathrm{\Psi }_{\frac{1}{2}}^{,,}(X)`$ is the class of pseudodifferential operators defined in §3.3.
### 5.2. Calculus in the model case
To define $`\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^𝔭(X)`$ we proceed locally and put $`\mathrm{\Sigma }`$ into a normal form $`\mathrm{\Sigma }_0=\{\xi _1=0\}`$ (locally).
The model case is obtained by taking symbols satisfying (5.3) globally and defining
$$a=a(x,\xi ,\lambda ,h,\stackrel{~}{h}),\lambda =\stackrel{~}{h}\xi _1/h,$$
satisfying
(5.5)
$$_{x_1}^k_x^{}^\alpha _\xi ^\beta _\lambda ^pa(x,\xi ,\lambda ;h)=𝒪(h^m\stackrel{~}{h}^{\stackrel{~}{m}})(\stackrel{~}{h}/h)^{(|\alpha |+|\beta |)/2}\lambda ^{k_2+k/2p/2},$$
which is the same as (5.1). We will write (5.5) as
$$a\stackrel{~}{𝒪}_{\frac{1}{2}}\left(h^m\stackrel{~}{h}^{\stackrel{~}{m}}\lambda ^{k_2}\right).$$
and define an exact quantization in the usual way,
(5.6)
$$\begin{array}{cc}& \stackrel{~}{\mathrm{Op}}_{h,\stackrel{~}{h}}(a)u(x)=\hfill \\ & \frac{1}{2\pi \stackrel{~}{h}}\frac{1}{(2\pi h)^{n1}}a(x,\xi ^{},(h/\stackrel{~}{h})\lambda ,\lambda ;h)e^{ix^{}y^{},\xi ^{}/h+ix_1y_1,\lambda /\stackrel{~}{h}}u(y)𝑑y𝑑\xi ^{}𝑑\lambda ,\hfill \end{array}$$
$`n=\text{dim}X`$, and where
$$\lambda =(\stackrel{~}{h}/h)\xi _1.$$
For $`a\stackrel{~}{𝒪}(\lambda ^{k_2})`$ and $`b\stackrel{~}{𝒪}(\lambda ^{k_2^{}})`$ we have
(5.9)
$$\begin{array}{c}\stackrel{~}{\mathrm{Op}}_{h,\stackrel{~}{h}}(a)\stackrel{~}{\mathrm{Op}}_{h,\stackrel{~}{h}}(b)=\stackrel{~}{\mathrm{Op}}_{h,\stackrel{~}{h}}(a\mathrm{}_{h,\stackrel{~}{h}}b),a\mathrm{}_{h,\stackrel{~}{h}}b=\stackrel{~}{𝒪}\left(\lambda ^{k_2+k_2^{}}\right),\\ a\mathrm{}_{h,\stackrel{~}{h}}b(x,\xi ,\lambda ;h)\underset{\alpha ^n}{}\frac{1}{\alpha !}(h_\xi ^{})^\alpha ^{}(h_{\xi _1}+\stackrel{~}{h}_\lambda )^{\alpha _1}aD_x^\alpha b,\end{array}$$
where the asymptotic sum is defined up to terms in
$$\stackrel{~}{𝒪}_{\frac{1}{2}}(\stackrel{~}{h}^{\mathrm{}}\lambda ^{k_2+k_2^{}}).$$
To see this we write $`\stackrel{~}{\mathrm{Op}}_{h,\stackrel{~}{h}}(a)`$ as a quantization of an operator valued symbol, $`\mathrm{Op}_h(a)(x_1,\lambda )`$,
$$\mathrm{Op}_h(a)(x_1,\lambda )=\frac{1}{(2\pi h)^{n1}}a(x,\xi ^{},(h/\stackrel{~}{h})\lambda ,\lambda ;h)e^{ix^{}y^{},\xi ^{}/h}𝑑\xi ^{},$$
so that
(5.10)
$$\stackrel{~}{\mathrm{Op}}_{h,\stackrel{~}{h}}(a)=\mathrm{Op}_h(a)(x_1,\stackrel{~}{h}D_{x_1}).$$
In view of (3.20) we have
$$a\stackrel{~}{𝒪}_{\frac{1}{2}}(\lambda ^{k_2})$$
$$_{x_1}^k_\lambda ^p\mathrm{Op}_h(^\alpha a)(x_1,\lambda )_{L^2(^{n1})L^2(^{n1})}C_{\alpha ,p,k}(\stackrel{~}{h}/h)^{|\alpha |/2}\lambda ^{k_2+k/2p/2}.$$
Observing that
$$\stackrel{~}{\mathrm{Op}}_{h,\stackrel{~}{h}}(a)\stackrel{~}{\mathrm{Op}}_{h,\stackrel{~}{h}}(b)=$$
$$(\mathrm{exp}(i\stackrel{~}{h}D_\lambda ,D_{y_1})\mathrm{Op}_ha(x_1,\lambda )\mathrm{Op}_hb(y_1,\lambda ^{})_{x_1=y_1,\lambda =\lambda ^{}})(x_1,\stackrel{~}{h}D_{x_1}),$$
proves (5.9).
For the operator $`\stackrel{~}{\mathrm{Op}}_{h,\stackrel{~}{h}}(a)`$, with $`a`$ supported in $`|\xi _1|C`$, we define its principal symbol as the equivalence class of $`a`$ in
$$\stackrel{~}{𝒪}_{\frac{1}{2}}(\lambda ^{k_2})/\stackrel{~}{𝒪}_{\frac{1}{2}}(\stackrel{~}{h}\lambda ^{k_2}).$$
In view of (5.9) this symbol map is a homomorphism onto the quotient of symbol spaces:
(5.11)
$$\stackrel{~}{\mathrm{Op}}_{h,\stackrel{~}{h}}(a)[a]\stackrel{~}{𝒪}_{\frac{1}{2}}(\lambda ^{k_2})/\stackrel{~}{𝒪}_{\frac{1}{2}}(\stackrel{~}{h}\lambda ^{k_2})$$
We note that in the local model we are not concerned here with the behaviour as $`|\xi |\mathrm{}`$.
It will be useful to have the analogue of Beals’s characterization of pseudodifferential operators by stability under taking commutators. It follows from the proof of the semiclassical analogue of Beals’s result and its adaptations in \[10, Chapter 8\] and \[32, Lemma 4.2\].
###### Proposition 5.2.
Let $`A=A_h:𝒮(^n)𝒮^{}(^n)`$ and put $`x^{}=(x_2,\mathrm{},x_n)`$. For $`p`$, we define the norms
$$u_{(p)}=\stackrel{~}{h}D_{x_1}^pu_{L^2}.$$
Then
$$\text{ }A=\stackrel{~}{\mathrm{Op}}_{h,\stackrel{~}{h}}(a)\text{ for }a=\stackrel{~}{𝒪}_{\frac{1}{2}}(\lambda ^k)\text{}a=a(x,\xi ^{},\lambda ;h,\stackrel{~}{h}),$$
if and only if for all $`N,p,q0`$ and every sequence $`\mathrm{}_1(x^{},\xi ^{}),\mathrm{},\mathrm{}_N(x^{},\xi ^{})`$ of linear forms on $`^{2(n1)}`$ there exists $`C>0`$ for which
(5.12)
$$\begin{array}{cc}& \mathrm{ad}_{\mathrm{}_1(x^{},hD_x^{})}\mathrm{}\mathrm{ad}_{\mathrm{}_N(x^{},hD_x^{})}(\mathrm{ad}_{\stackrel{~}{h}D_{x_1}})^p(\mathrm{ad}_{x_1})^qAu_{(q/2\mathrm{min}(k,0))}\hfill \\ & Ch^{N/2}\stackrel{~}{h}^{N/2+p+q}u_{p/2+(\mathrm{max}(k,0))}.\hfill \end{array}$$
###### Proof.
We first observe that for $`A=A=\stackrel{~}{\mathrm{Op}}_{h,\stackrel{~}{h}}(a)`$ for $`a=\stackrel{~}{𝒪}_{\frac{1}{2}}(\lambda ^k)`$, (5.12) follows from the calculus.
Let
$$V_hu(x_1,\stackrel{~}{x}^{})\stackrel{\mathrm{def}}{=}h^{\frac{n1}{4}}u(x_1,h^{\frac{1}{2}}\stackrel{~}{x}^{})$$
so that
$$A=V_h^1\stackrel{~}{A}V_h,$$
where, as in the proof of Proposition 3.5, we have
(5.13)
$$\begin{array}{cc}& \mathrm{ad}_{\mathrm{}_1(x^{},D_x^{})}\mathrm{}\mathrm{ad}_{\mathrm{}_N(x^{},D_x^{})}(\mathrm{ad}_{D_{x_1}})^p(\mathrm{ad}_{x_1})^q\stackrel{~}{A}u_{(q/2\mathrm{min}(k,0))}\hfill \\ & C\stackrel{~}{h}^{N/2+q}u_{p/2+(\mathrm{max}(k,0))}.\hfill \end{array}$$
We can write $`\stackrel{~}{A}=a_{h,\stackrel{~}{h}}(x,D_x)`$, so that, $`A=\stackrel{~}{\mathrm{Op}}_{h,\stackrel{~}{h}}(a)`$, where
$$a(x,\lambda ,\xi ^{};h)=a_{h,\stackrel{~}{h}}(x_1,h^{1/2}x^{};\stackrel{~}{h}^1\lambda ,h^{1/2}\xi ^{}).$$
The required estimate on $`a`$ becomes
(5.14)
$$_{x_1}^p_{\xi _1}^q_{\stackrel{~}{x}^{}}^\alpha ^{}_{\stackrel{~}{\xi }^{}}^\beta ^{}a_{h,h^{}}=𝒪(1)\stackrel{~}{h}^{(|\alpha ^{}|+|\beta ^{}|)/2+q}\stackrel{~}{h}\xi _1^{kq/2+p/2}$$
We have
(5.15)
$$\stackrel{~}{A}\psi ,\varphi =\frac{1}{(2\pi )^n}e^{ix,\xi }a_{h,\stackrel{~}{h}}(x,\xi )\widehat{\psi }(\xi )\overline{\varphi (x)}𝑑x𝑑\xi ,$$
with
$$\widehat{\psi }(\xi )=(\psi )(\xi )=e^{ix,\xi }\psi (x)𝑑x,$$
and for $`\varphi ,\psi 𝒮(^n)`$. Let us fix $`(x_0,\xi _0),(y_0,\eta _0)T^{}^n`$, and $`\lambda 1`$. With $`\varphi ,\psi 𝒮(^n)`$ we put
$$\begin{array}{cc}& \psi _{x_0,\xi _0}(x)=\lambda ^{\frac{1}{4}}\psi (\lambda ^{\frac{1}{2}}(x_1x_{0,1}),x^{}x_0^{})e^{ix,\xi _0},\hfill \\ & \varphi _{y_0,\eta _0}(x)=\lambda ^{\frac{1}{4}}\varphi (\lambda ^{\frac{1}{2}}(x_1y_{0,1}),x^{}y_0^{})e^{ixy_0,\eta _0}.\hfill \end{array}$$
We see that
$$\begin{array}{cc}& \widehat{\psi }_{x_0,\xi _0}(\xi )=\lambda ^{\frac{1}{4}}\widehat{\psi }(\lambda ^{\frac{1}{2}}(\xi _1\xi _{0,1}),\xi ^{}\xi _0^{})e^{i\xi \xi _0,x_0},\hfill \\ & \widehat{\varphi }_{y_0,\eta _0}(\xi )=\lambda ^{\frac{1}{4}}\widehat{\varphi }(\lambda ^{\frac{1}{2}}(\xi _1\eta _{0,1}),\xi ^{}\eta _0^{})e^{i\xi ,y_0}.\hfill \end{array}$$
We have
$$B\stackrel{\text{def}}{=}\text{ad}_x^{}^\alpha ^{}\text{ad}_{D_x^{}}^\beta ^{}\text{ad}_{x_1}^q\text{ad}_{D_{x_1}}^p\stackrel{~}{A}=\left(i\right)^{|\alpha ^{}|+|\beta ^{}|+q+p}b_{h,\stackrel{~}{h}}(x,D),$$
$$b_{h,\stackrel{~}{h}}(x,\xi )=(_\xi ^{})^\alpha ^{}_x^{}^\beta ^{}(_{\xi _1})^q_{x_1}^pa_{h,\stackrel{~}{h}}(x,\xi ).$$
Let us now assume we have the commutator estimate in the lemma with $`k0`$. Since
$$u_{(q/2)}=\stackrel{~}{h}D_{x_1}^{q/2}u_{L^2}$$
is the norm dual to $`_{(q/2)}`$, we get from (5.13)
(5.16)
$$|B\psi _{x_0,\xi _0},\varphi _{y_0,\eta _0}|\stackrel{~}{h}^{(|\alpha ^{}|+|\beta ^{}|)/2+q}\psi _{x_0,\xi _0}_{(p/2+k)}\varphi _{y_0,\eta _0}_{(q/2)}.$$
Let $`_{0,1}`$ denote the first component of $`_0^n`$. For fixed $`\psi ,\varphi 𝒮(^n)`$, we have
(5.19)
$$\begin{array}{c}\psi _{x_0,\xi _0}_{(p/2+k)}^2C_NI_{p/2+k}^N(\lambda ,\xi _{0,1}),\varphi _{y_0,\eta _0}_{(q/2)}^2C_NI_{q/2}^N(\lambda ,\eta _{0,1})\\ I_r^N(\lambda ,\tau )\stackrel{\mathrm{def}}{=}\lambda ^{\frac{1}{2}}_{}\stackrel{~}{h}\rho ^{2r}\lambda ^{\frac{1}{2}}(\rho \tau )^N𝑑\rho \end{array}$$
Using (5.15) we rewrite the left hand side of (5.16) as
$$\frac{1}{(2\pi )^n}\left|e^{ix,\xi }b_{h,\stackrel{~}{h}}(x,\xi )\widehat{\psi }_{x_0,\xi _0}(\xi )\overline{\varphi }_{y_0,\eta }(x)𝑑x𝑑\xi \right|.$$
Decomposing the first exponent in the integral as
$$x,\xi =y_0,\xi _0+xy_0,\xi _0+\xi \xi _0,y_0+xy_0,\xi \xi _0$$
and using the formulæ for $`\widehat{\psi }_{x_0,\xi _0}`$, $`\varphi _{y_0,\eta _0}`$, we rewrite it further as
$$\frac{1}{(2\pi )^n}|b_{h,\stackrel{~}{h}}(x,\xi )\mathrm{exp}(ixy_0,\xi \xi _0)\widehat{\psi }(\lambda ^{\frac{1}{2}}(\xi _1\xi _{0,1}),\xi ^{}\xi _0^{})$$
$$\overline{\varphi }(\lambda ^{\frac{1}{2}}(x_1y_{0,1}),x^{}y_0^{})\mathrm{exp}(i(\xi \xi _0,y_0x_0+xy_0,\xi _0\eta _0))dxd\xi |$$
Summing up, we get
(5.22)
$$\begin{array}{c}\left|\left(\chi b_{h,\stackrel{~}{h}}\right)(\eta _0\xi _0,x_0y_0)\right|=𝒪(1)\stackrel{~}{h}^{(|\alpha ^{}|+|\beta ^{}|)/2+q}\psi _{x_0,\xi _0}_{p/2+k}\varphi _{y_0,\eta _0}_{q/2},\\ \chi (x,\xi )=e^{ixy_0,\xi \xi _0}\widehat{\psi }(\lambda ^{\frac{1}{2}}(\xi _1\xi _{0,1}),\xi ^{}\xi _0^{})\overline{\varphi }(\lambda ^{\frac{1}{2}}(x_1y_{0,1}),x^{}y_0^{}).\end{array}$$
Writing
$$\zeta _1\stackrel{\mathrm{def}}{=}\frac{\eta _{0,1}\xi _{0,1}}{\lambda ^{\frac{1}{2}}},z_1\stackrel{\mathrm{def}}{=}\lambda ^{\frac{1}{2}}(x_{0,1}y_{0,1}),\zeta ^{}=\eta _0^{}\xi _0^{},z^{}=x_0^{}y_0^{},$$
we therefore get
$$\left(\stackrel{~}{\chi }\stackrel{~}{b}_{h,\stackrel{~}{h}}\right)(\zeta ,z)=𝒪(1)\stackrel{~}{h}^{(|\alpha ^{}|+|\beta ^{}|)/2+q}\psi _{x_0,\xi _0}_{(p/2+k)}\varphi _{y_0,\eta _0}_{(q/2)},$$
where
$$\stackrel{~}{b}_{h,\stackrel{~}{h}}(X,\mathrm{\Xi })\stackrel{\mathrm{def}}{=}b_{h,\stackrel{~}{h}}(y_0+\lambda ^{\frac{1}{2}}X,\xi _0+\lambda ^{\frac{1}{2}}\mathrm{\Xi }),$$
and
$$\stackrel{~}{\chi }(X,\mathrm{\Xi })\stackrel{\mathrm{def}}{=}e^{iX,\mathrm{\Xi }}\widehat{\psi }(\mathrm{\Xi })\overline{\varphi }(X).$$
We then have
$`\left(_\mathrm{\Xi }^{}^{\stackrel{~}{\alpha }^{}}_{\mathrm{\Xi }_1}^{\stackrel{~}{q}}_X^{}^{\stackrel{~}{\beta }^{}}_{X_1}^{\stackrel{~}{p}}\stackrel{~}{\chi }\stackrel{~}{b}_{h,\stackrel{~}{h}}\right)(\zeta ,z)=`$
$`𝒪_N(1)\stackrel{~}{h}^{(|\alpha ^{}|+|\beta ^{}|)/2+q}I_{(p/2+k+\stackrel{~}{p}/2)}^N(\lambda ,\xi _{0,1})I_{(q/2\stackrel{~}{q}/2)}(\lambda ,\eta _{0,1})\lambda ^{\stackrel{~}{p}/2+\stackrel{~}{q}/2},`$
which putting
$$\stackrel{~}{\alpha }=(\stackrel{~}{q},\stackrel{~}{\alpha }^{}),\stackrel{~}{\beta }=(\stackrel{~}{p},\stackrel{~}{\beta }^{}),$$
we rewrite as
$$\begin{array}{cc}& z^{\stackrel{~}{\alpha }}\zeta ^{\stackrel{~}{\beta }}\left(\stackrel{~}{\chi }\stackrel{~}{b}_{h,\stackrel{~}{h}}\right)(\zeta ,z)=\hfill \\ & 𝒪(1)\stackrel{~}{h}^{(|\alpha ^{}|+|\beta ^{}|)/2+q}I_{(p/2+k+\stackrel{~}{p}/2)}^N(\lambda ,\xi _{0,1})I_{(q/2\stackrel{~}{q}/2)}(\lambda ,\eta _{0,1})\lambda ^{\stackrel{~}{p}/2+\stackrel{~}{q}/2}\hfill \end{array}.$$
We now go back to (5.19) and choose $`\lambda =\stackrel{~}{h}\xi _{0,1}`$. Then
$$\begin{array}{cc}\hfill I_{(p/2+\stackrel{~}{p}/2+k)}^N(\stackrel{~}{h}\xi _{0,1},\xi _{0,1})& =\left(_{}\stackrel{~}{h}(\xi _{0,1}+\stackrel{~}{h}\xi _{0,1}^{\frac{1}{2}}r)^{p+\stackrel{~}{p}+2k}r^N𝑑r\right)^{\frac{1}{2}}\hfill \\ & =C_N(_{}R+R^{\frac{1}{2}}\stackrel{~}{h}r)^{p+\stackrel{~}{p}+2k}r^Ndr)^{\frac{1}{2}}_{R=\stackrel{~}{h}\xi _{0,1}}\hfill \\ & R^{(p+\stackrel{~}{p}+2k)/2}_{R=\stackrel{~}{h}\xi _{0,1}}=\stackrel{~}{h}\xi _{0,1}^{p/2+\stackrel{~}{p}/2+k}.\hfill \end{array}$$
It follows that for any $`\stackrel{~}{\alpha }`$ and $`N`$ we have
(5.23)
$$|z^{\stackrel{~}{\alpha }}\left(\stackrel{~}{\chi }\stackrel{~}{b}_{h,\stackrel{~}{h}}\right)(\zeta ,z)|=𝒪_N(1)\stackrel{~}{h}^{(|\alpha ^{}|+|\beta ^{}|)/2+q}I_{(q/2\stackrel{~}{q}/2)}^N\lambda ^{\stackrel{~}{q}/2+p/2+k}\zeta ^N,\lambda =\stackrel{~}{h}\xi _{0,1}.$$
We now integrate the left hand side in $`\zeta _1=(\eta _{0,1}\xi _{0,1})/\lambda ^{\frac{1}{2}}`$:
$$\begin{array}{cc}& _{}I_{q/2\stackrel{~}{q}/2}^N\frac{\eta _{0,1}\xi _{0,1}}{\stackrel{~}{h}\xi _{0,1}^{\frac{1}{2}}}^N\stackrel{~}{h}\xi _{0,1}^{\frac{1}{2}}𝑑\eta _{0,1}\hfill \\ & =_{}\left(_{}\stackrel{~}{h}(\eta _{0,1}+\stackrel{~}{h}\xi _{0,1}^{\frac{1}{2}}r)^{q\stackrel{~}{q}}r^N𝑑r\right)^{\frac{1}{2}}\frac{\eta _{0,1}\xi _{0,1}}{\stackrel{~}{h}\xi _{0,1}^N}^{\frac{1}{2}}\stackrel{~}{h}\xi _{0,1}^{\frac{1}{2}}𝑑\eta _{0,1}\hfill \\ & C_N^{}_{}\left(_{}\stackrel{~}{h}(\xi _{0,1}+\stackrel{~}{h}\xi _{0,1}^{\frac{1}{2}}(r+\zeta _1))^{q/2\stackrel{~}{q}/2}r^N𝑑r\right)^{\frac{1}{2}}\zeta _1^N𝑑\zeta _1\hfill \\ & \stackrel{~}{h}\xi _{0,1}^{q/2\stackrel{~}{q}/2}.\hfill \end{array}$$
Returning to (5.23) we see that
$$_^n|\left(\stackrel{~}{\chi }\stackrel{~}{b}_{h,\stackrel{~}{h}}\right)(\zeta ,z)|𝑑\zeta =𝒪(1)\stackrel{~}{h}^{(|\alpha ^{}|+|\beta ^{}|)/2+q}\stackrel{~}{h}\xi _{0,1}^{p/2+kq/2}z^N$$
and consequently
$$\stackrel{~}{b}_{h,\stackrel{~}{h}}\stackrel{~}{\chi }(0,0)=𝒪(1)\stackrel{~}{h}^{(|\alpha ^{}|+|\beta ^{}|)/2+q}\stackrel{~}{h}\xi _{0,1}^{k+p/2q/2}.$$
Combining this with the definition of $`b_{h,\stackrel{~}{h}}`$ and $`\stackrel{~}{b}_{h,\stackrel{~}{h}}`$ gives
$$_\xi ^{}^\alpha ^{}_{\xi _1}^q_x^{}^\beta ^{}_{x_1}^pa_{h,\stackrel{~}{h}}(y_0,\xi _0)=𝒪(1)\stackrel{~}{h}^{(|\alpha ^{}|+|\beta ^{}|)/2+q}\stackrel{~}{h}\xi _{0,1}^{k+p/2q/2},$$
which is (5.14) for $`k0`$. When $`k<0`$ we check that the assumption is satisfied for $`\stackrel{~}{h}D_{x_1}^kA`$ with $`k`$ replaced by $`0`$. The composition formula then gives the result. ∎
### 5.3. Modified Sobolev spaces
The norm used in Proposition 5.2 can be defined globally. We generalize it to include standard Sobolev spaces by adding additional information at a smooth compact hypersurface $`\mathrm{\Sigma }T^{}X`$.
Let $`Q\mathrm{\Psi }_h^{0,0}`$ be an operator, $`Q=Q^{}`$, with the principal symbol $`q`$ satisfying
(5.24)
$$\mathrm{\Sigma }=\{(x,\xi ):q(x,\xi )=0\},dq_\mathrm{\Sigma }0,|q(x,\xi )|1\text{ for }|\xi |C\text{ }.$$
For $`s,m`$ we define
(5.25)
$$H_\mathrm{\Sigma }^{s,m}(X)=\{u𝒮^{}(X):(\stackrel{~}{h}/h)Q^muH_h^s(X)\},u_{H_\mathrm{\Sigma }^{s,m}(X)}\stackrel{\mathrm{def}}{=}(\stackrel{~}{h}/h)Q^mu_{H_h^s(X)}.$$
Here $`H_h^s(X)`$ denotes the usual semiclassical Sobolev spaces defined in §3.1. The spaces are (complex) interpolation spaces, in $`m`$, and in $`s`$.
When $`m`$ the definition is equivalent to
$$H_\mathrm{\Sigma }^{s,m}(X)=\{\begin{array}{cc}\{u:(\stackrel{~}{h}/h)^kQ^kuH_h^s(X),0km\}\hfill & m0\hfill \\ & \\ \{u:u=_{k=0}^{|m|}(\stackrel{~}{h}/h)^kQ^ku_k,u_kH_h^s(X)\}\hfill & m0\hfill \end{array}$$
with the norms equivalent to the same natural way as for Sobolev spaces:
$$u_{H_\mathrm{\Sigma }^{s,m}(X)}\{\begin{array}{cc}_{k=0}^m(\stackrel{~}{h}/h)^kQ^ku_{H_h^s(X)},\hfill & m0\hfill \\ & \\ inf\{_{k=0}^{|m|}u_k_{H_h^s(X)}:u=_{k=0}^{|m|}(\stackrel{~}{h}/h)^kQ^ku_k\},\hfill & m0\hfill \end{array}$$
This can be seen using a spectral decomposition of $`Q`$ which is assumed to be bounded and self-adjoint. We use the following simplified notation
$$H_\mathrm{\Sigma }^m(X)\stackrel{\mathrm{def}}{=}H_\mathrm{\Sigma }^{0,m}(X),m.$$
The spaces $`H_\mathrm{\Sigma }^{s,m}(X)`$ have the following basic invariance property:
###### Lemma 5.3.
The definition (5.25) does not depend on the choice of $`Q\mathrm{\Psi }_h^{0,0}(X)`$ satisfying (5.24). If $`A\mathrm{\Psi }^{0,0}(X)`$ has the property that $`d(\mathrm{WF}_h(A),\mathrm{\Sigma })>1/C`$ then for $`M0`$
(5.26)
$$\begin{array}{cc}& Au_{H_\mathrm{\Sigma }^M(X)}(h/\stackrel{~}{h})^Mu_{L^2(X)},\hfill \\ & Au_{L^2(X)}(h/\stackrel{~}{h})^Mu_{H_\mathrm{\Sigma }^M(X)},\hfill \end{array}$$
for $`u𝒞_\mathrm{c}^{\mathrm{}}(X)`$.
Also, suppose that $`F`$ is a $`0`$th order $`h`$-Fourier Integral operator associated to a canonical transformation which maps $`\mathrm{\Sigma }`$ to another hypersurface $`\mathrm{\Sigma }^{}`$ satisfying our hypothesis. Then
$$F:H_\mathrm{\Sigma }^{s,m}(X)H_\mathrm{\Sigma }^{}^{s,m}(X).$$
###### Proof.
Let $`Q^{}`$ be another operator satisfying (5.24). Then
$$Q^{}=AQ+E=QA+E^{},$$
where $`A\mathrm{\Psi }^{0,0}(X)`$ is uniformly elliptic, and $`E,E^{}\mathrm{\Psi }^{1,1}(X)`$. Because of the interpolation property of $`H_\mathrm{\Sigma }^{s,m}`$ we only need to check the independence for $`k`$. Then, by induction on $`k`$,
$$(\stackrel{~}{h}/h)^k(Q^{})^ku_{H^s(X)}C_k\underset{\mathrm{}=0}{\overset{k}{}}\stackrel{~}{h}^k\mathrm{}(\stackrel{~}{h}/h)^{\mathrm{}}Q^{\mathrm{}}_{H^s(X)},$$
where $`C_k`$ depends only on $`A`$ and $`E`$. This shows that the definition for $`m0`$ is independent of the choice of $`Q`$. The case of $`m<0`$ is similar. To see the first inequality in (5.26) we recall that
$$Au_{H_\mathrm{\Sigma }^M(X)}=inf\{\underset{k=0}{\overset{M}{}}u_k_{L^2(X)}:Au=\underset{k=0}{\overset{M}{}}(\stackrel{~}{h}/h)^kQ^ku_k\}.$$
Because of the ellipticity of $`Q`$ on $`\mathrm{WF}_h(A)`$ we can find $`B`$ such that
$$Au=Q^kBAu+v,v_{L^2}=𝒪(h^{\mathrm{}})u_{L^2}.$$
Hence we can take $`u_0=v`$ and $`u_M=(h/\stackrel{~}{h})^MBAu`$, $`u_k=0`$ for all other $`k`$’s, so that
$$Au_{H_\mathrm{\Sigma }^M(X)}C(h/\stackrel{~}{h})^Mu_{L^2(X)}+𝒪(h^{\mathrm{}})u_{L^2(X)}$$
proving the first estimate in (5.26) ($`h<\stackrel{~}{h}`$ everywhere here). Since $`H_\mathrm{\Sigma }^M(X)`$ is the dual of $`H_\mathrm{\Sigma }^M`$ this is equivalent to
$$v𝒞_\mathrm{c}^{\mathrm{}}(X)Au,v_{L^2(X)}C(h/\stackrel{~}{h})^Mu_{L^2(X)}v_{H_\mathrm{\Sigma }^M(X)},$$
which in turn proves the second estimate with $`A`$ replaced by $`A^{}`$.
The Egorov theorem (Proposition 3.4 above) shows that $`QF`$ is equal to $`FQ^{}+E`$ where $`Q^{}`$ satisfies (5.24) with $`\mathrm{\Sigma }`$ replaced by $`\mathrm{\Sigma }^{}`$. The mapping property follows from the boundedness of $`F`$ on $`H^s(X)`$ and an argument similar to that above. ∎
### 5.4. Globally defined class of operators
For the global definition we cannot use the classes
$$\stackrel{~}{\mathrm{Op}}_{h,\stackrel{~}{h}}\left(\stackrel{~}{𝒪}_{\frac{1}{2}}(\lambda ^m)\right)$$
due to the presence of propagation in the $`h`$-sense. To see the problem let us consider the operator $`P`$ introduced in §1 and define
$$A=h^m\chi \left(\frac{\stackrel{~}{h}P}{h}\right),\chi 𝒞_\mathrm{c}^{\mathrm{}}().$$
In the model theory of §5.2 this operator is obtained from the symbol $`h^m\chi (\lambda )`$. However, when $`\stackrel{~}{h}`$ is fixed, the operator $`A`$ is an $`h`$-Fourier integral operator:
$$A=\frac{h^m}{2\pi \stackrel{~}{h}}_{}\widehat{\chi }(t/\stackrel{~}{h})\mathrm{exp}(itP/h)𝑑t,$$
associated to the relation
$$\{(\rho ,\rho ^{}):t,\mathrm{exp}(tH_p)(\rho )=\rho ^{},p(\rho )=p(\rho ^{})=0\}.$$
Hence the presence of almost closed orbits of the flow prevents a definition which would be purely local in the $`𝒪(h^{\mathrm{}})`$ sense. We observe however that the non-local contributions in $`A`$ are of the order $`𝒪(\stackrel{~}{h}^{\mathrm{}})`$.
These global concerns suggest the following definition of the residual class. First we introduce a useful cut-off operator. Let $`V_1,V_2`$ be two open neighbourhoods of $`\mathrm{\Sigma }`$, satisfying
$$\overline{V_1}V_2.$$
Then we choose
(5.27)
$$\gamma _\mathrm{\Sigma }\mathrm{\Psi }^{0,0}(X),\mathrm{WF}_h(\gamma _\mathrm{\Sigma })V_2,\mathrm{WF}_h(I\gamma _\mathrm{\Sigma })\mathrm{}V_1.$$
We now define the spaces of operators. Let
$$𝔭=(m,\stackrel{~}{m},k_1,k_2),𝔭_{\mathrm{}}=(m,\mathrm{},\mathrm{},k_2).$$
As before we start with the definition of a residual class:
Definition 1. We say that $`A\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^𝔭_{\mathrm{}}(X)`$ if
$$A:𝒞_\mathrm{c}^{\mathrm{}}(X)𝒞^{\mathrm{}}(X),(I\gamma _\mathrm{\Sigma })A,A(I\gamma _\mathrm{\Sigma })\mathrm{\Psi }_{\frac{1}{2}}^{m+k_2,\mathrm{},\mathrm{}}(X),$$
and for any $`u𝒞_\mathrm{c}^{\mathrm{}}(X)`$, any sequence $`\{b_j\}_{j=1}^NS^0(T^{}X)`$, and any $`k`$, $`p`$, and $`M`$, we have
$$\mathrm{ad}_{\mathrm{Op}_h^w(b_1)}\mathrm{}\mathrm{ad}_{\mathrm{Op}_h^w(b_N)}\mathrm{ad}_Q^kAu_{H_\mathrm{\Sigma }^{p+N/2k_2}(X)}Ch^{m+k}\stackrel{~}{h}^Mu_{H_\mathrm{\Sigma }^{p+k/2}},$$
where $`Q`$ is the operator in (5.25).
Definition 2. We say that $`A\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^𝔭(X)`$ if
$$A:𝒞_\mathrm{c}^{\mathrm{}}(X)𝒞^{\mathrm{}}(X),(I\gamma _\mathrm{\Sigma })A,A(I\gamma _\mathrm{\Sigma })\mathrm{\Psi }_{\frac{1}{2}}^{m+k_2,\stackrel{~}{m}k_2,k_1}(X),$$
for some cut-off operator $`\gamma _\mathrm{\Sigma }`$ satisfying (5.27).
* For any $`\chi 𝒞_\mathrm{c}^{\mathrm{}}(T^{}X)`$,
$$A=A_\chi ^{\mathrm{}}+A_\chi ^{\mathrm{}},A_\chi ^{\mathrm{}}\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^𝔭_{\mathrm{}}(X),$$
so that for any $`p`$, and any sequences
$$\{b_j\}_{j=1}^N,\{a_j\}_{j=1}^MS^0(T^{}X),H_pa_j_{\mathrm{supp}\chi }0,$$
$$\mathrm{ad}_{\mathrm{Op}_h^w(b_1)}\mathrm{}\mathrm{ad}_{\mathrm{Op}_h^w(b_N)}\mathrm{ad}_{\mathrm{Op}_h^w(a_1)}\mathrm{}\mathrm{ad}_{\mathrm{Op}_h^w(a_M)}\mathrm{ad}_Q^k\chi ^w\gamma _\mathrm{\Sigma }A_\chi ^{\mathrm{}}\chi ^wu_{H_\mathrm{\Sigma }^{p+N/2k_2}(X)}$$
$$Ch^{M/2+km}\stackrel{~}{h}^{M/2+N\stackrel{~}{m}}u_{H_\mathrm{\Sigma }^{p+k/2}(X)},\chi ^w=\mathrm{Op}_h^w(\chi ),$$
for any $`u𝒞_\mathrm{c}^{\mathrm{}}(X)`$.
* For any $`\psi _1,\psi _2𝒞_\mathrm{b}^{\mathrm{}}(T^{}X)`$, $`\mathrm{supp}\psi _1\mathrm{supp}\psi _2=\mathrm{}`$,
$$\mathrm{Op}_h^w(\psi _1)A\mathrm{Op}_h^w(\psi _2)\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^𝔭_{\mathrm{}}(X).$$
It is important to record
###### Lemma 5.4.
Definitions 1 and 2 are independent of the choice of the cut-off operator satisfying (5.27).
###### Proof.
Suppose that $`\gamma _\mathrm{\Sigma }^{}`$ is another cut-off operator satisfying (5.27). We need to show that
(5.28)
$$(\gamma _\mathrm{\Sigma }\gamma _\mathrm{\Sigma }^{})A,A(\gamma _\mathrm{\Sigma }\gamma _\mathrm{\Sigma }^{})\mathrm{\Psi }_{\frac{1}{2}}^{m+k_2,\stackrel{~}{m}k_2,\mathrm{}}(X),$$
and for that we will use only the commutators involving $`b_j`$’s in Definitions 1 and 2. The first inclusion is, in view of Lemma 3.5 and the fact that $`\gamma _\mathrm{\Sigma }\gamma _\mathrm{\Sigma }^{}`$ is smoothing, equivalent to
(5.31)
$$\begin{array}{c}\mathrm{ad}_{\mathrm{Op}_h^w(b_1)}\mathrm{}\mathrm{ad}_{\mathrm{Op}_h^w(b_N)}(\gamma _\mathrm{\Sigma }\gamma _\mathrm{\Sigma }^{})Au_{L^2(X)}\\ Ch^{mk_2+N/2}\stackrel{~}{h}^{\stackrel{~}{m}+k_2+N/2}u_{L^2(X)}.\end{array}$$
To see this we first use (5.26) to see that if $`E\mathrm{\Psi }^{0,0}(X)`$ and $`d(\mathrm{WF}_h(E),\mathrm{\Sigma })>1/C`$ then for any subsequence $`\{i_j\}_{j=1}^J`$ of $`1,\mathrm{},N`$,
(5.32)
$$\begin{array}{cc}& E\mathrm{ad}_{\mathrm{Op}_h^w(b_{i_1})}\mathrm{}\mathrm{ad}_{\mathrm{Op}_h^w(b_{i_J})}Au_{L^2(X)}\hfill \\ & C(h/\stackrel{~}{h})^{J/2k_2}\mathrm{ad}_{\mathrm{Op}_h^w(b_{i_1})}\mathrm{}\mathrm{ad}_{\mathrm{Op}_h^w(b_{i_J})}Au_{H_\mathrm{\Sigma }^{J/2k_2}(X)}\hfill \\ & Ch^m(h/\stackrel{~}{h})^{J/2k_2}\stackrel{~}{h}^{J\stackrel{~}{m}}u_{L^2(X)}=Ch^{mk_2+J/2}\stackrel{~}{h}^{\stackrel{~}{m}+k_2+J/2}u_{L^2(X)}.\hfill \end{array}$$
Using the derivation property of $`\mathrm{ad}_{\mathrm{Op}_h^w(b_j)}`$ and the fact that
$$\mathrm{ad}_{\mathrm{Op}_h^w(b_j)}(\gamma _\mathrm{\Sigma }\gamma _\mathrm{\Sigma }^{})=hE_j,E_j\mathrm{\Psi }^{0,0}(X),d(\mathrm{WF}_h(E_j),\mathrm{\Sigma })>1/C,$$
we can estimate the left hand side of (5.31) by a linear combination of terms of the form
$$h^{NJ}E\mathrm{ad}_{\mathrm{Op}_h^w(b_{i_1})}\mathrm{}\mathrm{ad}_{\mathrm{Op}_h^w(b_{i_J})}Au_{L^2(X)},E\mathrm{\Psi }^{0,0}(X),d(\mathrm{WF}_h(E),\mathrm{\Sigma })>1/C.$$
Consequently (5.31) follows from (5.32).
Since, by the calculus of §3.3
$$B\mathrm{\Psi }_{\frac{1}{2}}^{m+k_2,\mathrm{},\mathrm{}}(X)B^{}\mathrm{\Psi }_{\frac{1}{2}}^{m+k_2,\mathrm{},\mathrm{}}(X),$$
and since
$$\mathrm{ad}_{\mathrm{Op}_h^w(b_1)}\mathrm{}\mathrm{ad}_{\mathrm{Op}_h^w(b_N)}B^{}=(1)^N\left(\mathrm{ad}_{\mathrm{Op}_h^w(\overline{b}_1)}\mathrm{}\mathrm{ad}_{\mathrm{Op}_h^w(\overline{b}_N)}B\right)^{},$$
the second inclusion in (5.28) follows from the first one. ∎
We now have a natural mapping and invariance properties:
###### Proposition 5.5.
The operators $`\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^{0,0,0,0}(X)`$ form an algebra and $`\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^{0,\mathrm{},\mathrm{},0}(X)`$ is an ideal in that algebra. If $`A\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^𝔭_{\mathrm{}}(X)`$ then
$$A=𝒪(h^m\stackrel{~}{h}^M):H_\mathrm{\Sigma }^{M,r}(X)H_\mathrm{\Sigma }^{M,rk_2}(X),$$
for any $`m,M`$. If $`A\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^𝔭(X)`$ then
$$A=𝒪(h^m\stackrel{~}{h}^{\stackrel{~}{m}}):H_\mathrm{\Sigma }^{s+k_1,p+k_2}(X)H_\mathrm{\Sigma }^{s,p}(X),$$
for any $`s`$ and $`p`$.
Also, suppose that $`F,G`$ are $`0`$th order $`h`$-Fourier Integral operators, $`F`$ associated to a canonical transformation which maps $`\mathrm{\Sigma }`$ to another hypersurface $`\mathrm{\Sigma }^{}`$ satisfying our hypothesis, and $`G`$ to its inverse (both transformation need to be defined only locally). Then
(5.33)
$$A\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^𝔭(X)FAG\mathrm{\Psi }_{\mathrm{\Sigma }^{},\frac{1}{2}}^𝔭(X).$$
###### Proof.
We first show that $`\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^{0,\mathrm{},\mathrm{},0}(X)`$ is an algebra. In fact, in the notation of Definition 1,
$$\mathrm{ad}_{\mathrm{Op}_h^w(b_1)}\mathrm{}\mathrm{ad}_{\mathrm{Op}_h^w(b_N)}\mathrm{ad}_Q^kABu_{H_\mathrm{\Sigma }^{p+N/2}(X)},A,B\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^{0,\mathrm{},\mathrm{},0}(X),$$
is bounded by a linear combination of
$$\mathrm{ad}_{\mathrm{Op}_h^w(b_{i_1})}\mathrm{}\mathrm{ad}_{\mathrm{Op}_h^w(b_{i_J})}\mathrm{ad}_Q^lA\mathrm{ad}_{\mathrm{Op}_h^w(b_{k_1})}\mathrm{}\mathrm{ad}_{\mathrm{Op}_h^w(b_{i_{NJ}})}\mathrm{ad}_Q^{kl}Bu_{H_\mathrm{\Sigma }^{p+N/2}(X)},$$
$$\{i_j\}_{j=1}^J\{k_j\}_{j=1}^{NJ}=\{1,\mathrm{},N\}.$$
These terms are estimated by
$$Ch^l\stackrel{~}{h}^{M_1}\mathrm{ad}_{\mathrm{Op}_h^w(b_{k_1})}\mathrm{}\mathrm{ad}_{\mathrm{Op}_h^w(b_{i_{NJ}})}\mathrm{ad}_Q^{kl}Bu_{H_\mathrm{\Sigma }^{p+(NJ)/2+l/2}(X)}$$
$$C^{}h^k\stackrel{~}{h}^{M_1+M_2}u_{p+N/2+k},$$
where we used Definition 1 with $`A`$ and with $`B`$. Here $`M_1`$ and $`M_2`$ are any large integers. Checking that
$$(I\gamma _\mathrm{\Sigma })AB,AB(I\gamma _\mathrm{\Sigma })\mathrm{\Psi }_{\frac{1}{2}}^{0,\mathrm{},\mathrm{}}(X),$$
is the same as in the proof of Lemma 5.4.
We now check that for $`B\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^{0,\mathrm{},\mathrm{},0}(X)`$ and $`A\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^{0,0,0,0}(X)`$ we have
$$BA\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^{0,\mathrm{},\mathrm{},0}(X).$$
In the notation of Definition 2, we can write
$$B=\gamma _\mathrm{\Sigma }B_1\gamma _\mathrm{\Sigma }+(1\gamma _\mathrm{\Sigma })B_1+\gamma _\mathrm{\Sigma }B_1(1\gamma _\mathrm{\Sigma })+B_2,$$
and it suffices to check that
(5.34)
$$A\gamma _\mathrm{\Sigma }B_1\gamma _\mathrm{\Sigma }\mathrm{\Psi }_{\frac{1}{2}}^{0,\mathrm{},\mathrm{}}(X).$$
Since $`\mathrm{\Sigma }`$ is compact (and hence we can take $`\chi ^w=\gamma _\mathrm{\Sigma }`$) and Definition 2 is independent of the choice of $`\gamma _\mathrm{\Sigma }`$ we have, as its special case,
$$\mathrm{ad}_{\mathrm{Op}_h^w(b_1)}\mathrm{}\mathrm{ad}_{\mathrm{Op}_h^w(b_N)}\mathrm{ad}_Q^k\gamma _\mathrm{\Sigma }B_1\gamma _\mathrm{\Sigma }u_{H_\mathrm{\Sigma }^{p+N/2}(X)}Ch^k\stackrel{~}{h}^Nu_{H_\mathrm{\Sigma }^{p+k/2}(X)},$$
and the verification of (5.34) follows the proof of the algebra property of $`\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^{0,\mathrm{},\mathrm{},0}(X)`$.
To conclude the algebraic part of the proof we show that $`\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^{0,0,0,0}(X)`$ is closed under composition of operators. Let $`B_1,B_2\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^{0,0,0,0}(X)`$ and let $`\chi 𝒞_\mathrm{c}^{\mathrm{}}(T^{}X)`$. Since we already established that composition with elements of $`\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^{0,\mathrm{},\mathrm{},0}(X)`$ produces an operator in that space we only have to show that
$$B_1B_2=A_1+A_2,A_2\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^{0,\mathrm{},\mathrm{},0}(X),$$
so that for any $`p`$, and any sequences
$$\{b_j\}_{j=1}^N,\{a_j\}_{j=1}^MS^0(T^{}X),H_pa_j_{\mathrm{supp}\chi }0,$$
and
$$\mathrm{ad}_{\mathrm{Op}_h^w(b_1)}\mathrm{}\mathrm{ad}_{\mathrm{Op}_h^w(b_N)}\mathrm{ad}_{\mathrm{Op}_h^w(a_1)}\mathrm{}\mathrm{ad}_{\mathrm{Op}_h^w(a_M)}\mathrm{ad}_Q^k\chi ^w\gamma _\mathrm{\Sigma }A_1\chi ^wu_{H_\mathrm{\Sigma }^{p+N/2k_2}(X)}$$
$$Ch^{M/2+km}\stackrel{~}{h}^{M/2+N\stackrel{~}{m}}u_{H_\mathrm{\Sigma }^{p+k/2}(X)}$$
To find the decomposition of $`B_1B_2`$ we introduce $`\chi _j𝒞_\mathrm{c}^{\mathrm{}}(T^{}X)`$ such that $`\chi _11`$ on $`\mathrm{supp}\chi _0`$, and $`\chi _01`$ on $`\mathrm{supp}\chi `$. We choose $`\mathrm{supp}\chi _1`$ sufficiently close to the support of $`\chi `$ so that the functions $`a_j`$ satisfy $`H_pa_j_{\mathrm{supp}\chi _1}0`$. We then put
$$A_2=\chi _0^w(B_1)_{\chi _1}^{\mathrm{}}(1\chi _1^w)(B_2)_{\chi _1}^{\mathrm{}}+B_1(B_2)_{\chi _1}^{\mathrm{}}+(B_1)_{\chi _1}^{\mathrm{}}B_2,$$
which is in $`\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^{0,\mathrm{},\mathrm{},0}(X)`$ since $`\mathrm{supp}\chi _0\mathrm{supp}(1\chi _1^w)=\mathrm{}`$, and we can use the second part of Definition 2.
We have
$$A_1=(1\chi _0^w)(B_1)_{\chi _1}^{\mathrm{}}(B_2)_{\chi _1}^{\mathrm{}}+\chi _0^w(B_1)_{\chi _1}^{\mathrm{}}\chi _1^w(B_1)_{\chi _1}^{\mathrm{}}.$$
Up to negligible, $`𝒪(h^{\mathrm{}})`$, errors
$$\chi ^w(1\chi _0^w)0,\chi ^w\chi _0^w\chi ^w.$$
Hence we need to check that
$$\mathrm{ad}_{\mathrm{Op}_h^w(b_1)}\mathrm{}\mathrm{ad}_{\mathrm{Op}_h^w(b_N)}\mathrm{ad}_{\mathrm{Op}_h^w(a_1)}\mathrm{}\mathrm{ad}_{\mathrm{Op}_h^w(a_M)}\mathrm{ad}_Q^k\chi ^w(B_1)_{\chi _1}^{\mathrm{}}\chi _0^w(B_2)_{\chi _1}^{\mathrm{}}\chi ^w$$
$$=𝒪(h^{M/2+km}\stackrel{~}{h}^{M/2+N\stackrel{~}{m}}):H_\mathrm{\Sigma }^{p+N/2k_2}(X)H_\mathrm{\Sigma }^{p+k/2}(X),$$
and this follows from the Leibnitz rule for $`\mathrm{ad}_{b_j}`$ and assumptions on $`(B_j)_{\chi _1}^{\mathrm{}}`$.
The mapping properties are immediate from the definitions: we apply them with no commutators, and in the case of Definition 2, with $`\chi ^w=\gamma _\mathrm{\Sigma }`$.
Lemma 5.3 shows that the spaces $`H_\mathrm{\Sigma }^{\pm M}`$ transform correctly under $`F`$ and hence Proposition 3.4 (Egorov’s Theorem) shows that for
$$A\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^𝔭_{\mathrm{}}(X)FAG\mathrm{\Psi }_{\mathrm{\Sigma }^{},\frac{1}{2}}^𝔭_{\mathrm{}}(X),$$
that is (5.33) holds for the residual class. Since the conditions on $`a_j`$, $`b_j`$’s in Definition 2 are symplectically invariant, Egorov’s theorem (Proposition 3.4) again shows that $`FAG\mathrm{\Psi }_\mathrm{\Sigma }^{}^𝔭(X)`$. ∎
### 5.5. The symbol map
To define the symbol map we will use the invariance given by Proposition 5.5 and the symbol map in the model case. We start with
###### Lemma 5.6.
In the notation of Definition 2, suppose that $`X=T^{}^n`$ and that
$$\mathrm{\Sigma }V=\{\xi _1=0\}V,VT^{}^n,\text{ is open, and }\mathrm{supp}\chi V.$$
Then
$$\mathrm{Op}_h^w(\chi )A_1\mathrm{Op}_h^w(\chi )=\stackrel{~}{\mathrm{Op}}_{h,\stackrel{~}{h}}(a_\chi ),a_\chi =\stackrel{~}{𝒪}_{\frac{1}{2}}(h^m\stackrel{~}{h}^{\stackrel{~}{m}}\lambda ^{k_2}).$$
###### Proof.
This is a consequence of Proposition 5.2 and the properties of the term $`A_1`$ in Definition 2. ∎
To construct a symbol map, that is a homomorphism
$$\sigma _{\mathrm{\Sigma },h,\stackrel{~}{h}}:\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^𝔭(X)S_{\mathrm{\Sigma },\frac{1}{2}}^𝔭(T^{}X)/S_{\mathrm{\Sigma },\frac{1}{2}}^𝔭^{}(T^{}X),$$
$$𝔭=(m,\stackrel{~}{m},k_1,k_2),𝔭^{}=(m,\stackrel{~}{m}1,k_11,k_2),$$
such that the sequence
$$0\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^𝔭^{}(X)\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^𝔭(X)\stackrel{\sigma _{\mathrm{\Sigma },h,\stackrel{~}{h}}}{}S_{\mathrm{\Sigma },\frac{1}{2}}^𝔭(T^{}X)/S_{\mathrm{\Sigma },\frac{1}{2}}^𝔭^{}(T^{}X)0$$
is exact, for an arbitrary $`\mathrm{\Sigma }`$ we will use Lemma 5.6. That requires putting $`\mathrm{\Sigma }`$ locally to the model hypersurface $`\xi _1=`$ which on the quantum level is done using $`h`$-Fourier integral operators. Hence we need a local invariance statement given in the next
###### Proposition 5.7.
Let $`U`$ be an $`h`$-Fourier Integral Operator, elliptic in $`V\times V`$, where $`V`$ is a neighbourhood of $`(0,0)T^{}^n`$, the compact hypersurface $`\mathrm{\Sigma }`$ satisfies,
$$\mathrm{\Sigma }V=\{\xi _1=0\}V.$$
Assume also that the canonical transformation associated to $`U`$, $`\kappa `$, satisfies:
(5.35)
$$\kappa (0,0)=(0,0),\kappa (\{\xi _1=0\}V)\{\xi _1=0\}.$$
Let $`A=\stackrel{~}{\mathrm{Op}}_{h,\stackrel{~}{h}}(a)`$ where $`a=\stackrel{~}{𝒪}_{\frac{1}{2}}(h^m\stackrel{~}{h}^{\stackrel{~}{m}}\lambda ^k),a(x,\xi ,\lambda )0\text{ for }(x,\xi )V\text{.}`$
If $`U^1`$ is the microlocal inverse of $`U`$ near $`V\times V`$, then
$$U^1AU=\stackrel{~}{\mathrm{Op}}_{h,\stackrel{~}{h}}(b)+E,b=aK,$$
$$E\mathrm{\Psi }_\mathrm{\Sigma }^𝔭^{}(^n),𝔭^{}=(m,\stackrel{~}{m}1,\mathrm{},k1),$$
where $`K`$ is the natural lift of $`\kappa `$ to the $`(x,\xi ,\lambda )`$ variables:
(5.36)
$$K(y,\eta ,\mu )=(x,\xi ,\lambda )(x,\xi )=\kappa (y,\eta ),\lambda =(\xi _1/\eta _1)\mu .$$
###### Proof.
We start by observing that the proposition holds in the special cases $`\kappa (x,\xi )=(x,\xi )`$ and $`\kappa (x,\xi )=(x,\xi )`$. The first special case concerns conjugation with elliptic classical $`h`$-pseudodifferential operators and it follows from the discussion after (5.9). The second special case follows from the first one and the fact that the proposition is easily checked for $`Uu(x)=u(x)`$. As a consequence we can asssume that $`\kappa `$ preserves the sign of $`\xi _1`$:
(5.37)
$$\kappa (y,\eta )=(x,\xi )\xi _1\eta _10.$$
We will prove the proposition by a deformation method inspired by the “Heisenberg picture of quantum mechanics” and for that we need the following geometric
###### Lemma 5.8.
Let $`\kappa `$ be a smooth canonical transformation satisfying (5.35) and (5.37). Then we can find a piecewise smooth family of canonical transformations $`[0,1]t\kappa _t`$ satisfying (5.35), (5.37) and such that $`\kappa _0=id`$ and $`\kappa _1=\kappa `$.
###### Proof.
Let us denote by $`\mathrm{\Sigma }`$ the hypersurface given by $`\xi _1=0`$. We first observe that if $`\kappa `$ is a linear symplectic transformation preserving $`\mathrm{\Sigma }`$ then we can find a family of linear symplectic transformations, $`\kappa _t`$, satisfying the conclusions of the lemma: the subgroup of elements of $`Sp(n,)`$ preserving a half space bounded by $`\mathrm{\Sigma }`$ is connected.
Hence we can assume that $`d\kappa (0,0)=Id`$. Now introduce
$$\kappa _t(x,\xi )=t^1\kappa (tx,t\xi ),$$
a smooth family of symplectic transformations, preserving a half space bounded by $`\mathrm{\Sigma }`$, with $`\kappa _1=\kappa `$ and $`\kappa _0=Id`$. ∎
We now return to the proof of Proposition 5.7. For simplicity we can assume that $`m=\stackrel{~}{m}=0`$. Let $`\kappa _t`$ be a piecewise smooth family with $`\kappa _1=\kappa `$, $`\kappa _0=id`$. Using (3.8) we construct a piecewise smooth family of classical elliptic $`h`$-Fourier integral operators, $`U_t`$, defined microlocally near $`(0,0)`$ and associated to $`\kappa _t`$. If we demand that $`U_1=U`$ then $`U_0`$ is a pseudodifferential operator elliptic at $`(0,0)`$. For notational convenience we assume that our deformation is smooth in $`t`$ – the piecewise smooth case follows from the same argument applied in several steps and that $`U_0=Id`$ (the last condition can be arranged by composing $`U`$ with a elliptic pseudodifferential operator). Thus we have
(5.38)
$$hD_tU_t+U_tQ_t=0,$$
where $`Q_t`$ is a smooth family of classical $`h`$-pseudodifferential operators of order $`0`$ with the leading symbol $`q_t`$ satisfying
(5.39)
$$\frac{d}{dt}\kappa _t(x,\xi )=(\kappa _t)_{}\left(H_{q_t}(x,\xi )\right),$$
in a neighbourhood of $`V`$. It follows from (5.39) that $`H_{q_t}`$ is tangent to $`\mathrm{\Sigma }`$ and hence
(5.40)
$$_{x_1}q_t(x,\xi )=\xi _1r_t(x,\xi ).$$
We extend $`q_t`$ to a globally defined function in $`S(T^{}^n,x,\xi )`$, keeping the property (5.40). This defines a family of global canonical transformation which coincide with $`\kappa _t`$ near $`V`$.
Let $`V_t`$ satisfy
(5.41)
$$hD_tV_t=Q_tV_t,V_0=Id.$$
It follows that $`V_t=U_t^1`$, and if we take $`Q_t`$ to be self-adjoint, $`V_t^{}=U_t`$.
We will now use (5.38) and (5.39) to prove Egorov’s theorem (3.10) for the new class. Thus we consider $`A_t=U_t^1AU_t`$, so that $`A_1`$ is the operator we want to study and $`A_0=A`$ is the given operator. From Proposition 5.5 we already know that $`A_t\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^{0,0,0,k}(^n)`$. Using Lemma 5.6 we can write
$$A_t\stackrel{~}{\mathrm{Op}}_{h,\stackrel{~}{h}}(a_t)mod\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^{0,\mathrm{},\mathrm{},k}(^n),$$
microlocally near $`V`$. From (5.38) and (5.41), we get
(5.45)
$$\begin{array}{c}\{\begin{array}{c}hD_tA_t=[Q_t,A_t],\hfill \\ A_0=A,\hfill \end{array}\end{array}$$
To compute the commutator on the symbolic level we need:
###### Lemma 5.9.
Suppose that $`a\stackrel{~}{𝒪}_{\frac{1}{2}}(\lambda ^k)`$ and that $`bS^{0,\mathrm{}}(T^{}^n)`$ (that is, $`b`$ is a symbol in the sense of §3.1), does not depend on $`\lambda `$. In addition let us assume that
$$_{x_1}b=\xi _1r(x,\xi ).$$
Then
(5.46)
$$\frac{i}{h}[\stackrel{~}{\mathrm{Op}}_{h,\stackrel{~}{h}}(b),\stackrel{~}{\mathrm{Op}}_{h,\stackrel{~}{h}}(a)]=\stackrel{~}{\mathrm{Op}}_{h,\stackrel{~}{h}}(c),c=(H_br\lambda _\lambda )a+\stackrel{~}{𝒪}_{\frac{1}{2}}(\stackrel{~}{h}\lambda ^k).$$
###### Proof.
We will use (5.10) to compute $`a\mathrm{}_{h,\stackrel{~}{h}}b`$ (and $`b\mathrm{}_{h,\stackrel{~}{h}}a`$), noting that
$$\mathrm{Op}_ha(x_1,\stackrel{~}{h}D_{x_1})\mathrm{Op}_hb(x_1,\stackrel{~}{h}D_{x_1})=\left(\mathrm{Op}_h(a)\mathrm{}_{\stackrel{~}{h}}\mathrm{Op}_h(b)\right)(x_1,\stackrel{~}{h}D_{x_1}).$$
Since $`bS^{0,\mathrm{}}`$, Lemma 3.6 shows that the composition formula in the $`(x^{},\xi ^{})`$ is given by an asymptotic series in $`(h\stackrel{~}{h})^{\frac{1}{2}}`$, and a $`𝒪(h^{\mathrm{}})`$ error. The only subtlety lies in the dependence on $`(x_1,\lambda )`$ and to explain it we suppress the other variables. We have
$$_{x_1}^{\mathrm{}}_\lambda ^pb(x_1,\lambda )=\{\begin{array}{cc}𝒪((h/\stackrel{~}{h})^p(h/\stackrel{~}{h})\lambda ^{\mathrm{}})\hfill & \mathrm{}=0\hfill \\ 𝒪((h/\stackrel{~}{h})^{p+1}\lambda (h/\stackrel{~}{h})\lambda ^{\mathrm{}})\hfill & \mathrm{}>0,\hfill \end{array}$$
and $`_{x_1}^p_\lambda ^{\mathrm{}}a=𝒪(\lambda ^{k\mathrm{}/2+p/2})`$. Consequently the terms in the expansions of $`a\mathrm{}_{\stackrel{~}{h}}b`$ and $`b\mathrm{}_{\stackrel{~}{h}}a`$ are bounded by
$$C_ph\stackrel{~}{h}^{p1}\lambda ^{kp/2+1}(h/\stackrel{~}{h})\lambda ^{\mathrm{}},C_p(h/\stackrel{~}{h})^p\stackrel{~}{h}^p\lambda ^{k+p/2}(h/\stackrel{~}{h})\lambda ^{\mathrm{}},p>0,$$
respectively. Hence we have expansions such that for $`p>1`$ the terms are bounded by $`h\stackrel{~}{h}^{p1}\lambda ^k`$, and the errors are $`𝒪(h\stackrel{~}{h}^{\mathrm{}}\lambda ^k)`$. This argument shows that
$$a\mathrm{}_hb=ab+h\left(\underset{i=1}{\overset{n}{}}_{\xi _i}aD_{x_i}b+r\lambda _\lambda a\right)+\stackrel{~}{𝒪}_{\frac{1}{2}}\left(h\stackrel{~}{h}\lambda ^k\right),$$
and
$$b\mathrm{}_ha=ab+h\underset{i=1}{\overset{n}{}}_{\xi _i}bD_{x_i}a+\stackrel{~}{𝒪}_{\frac{1}{2}}\left(h\stackrel{~}{h}\lambda ^k\right),$$
from which the lemma follows. ∎
Since $`q_t`$ satisfies (5.40), Lemma 5.9 gives
$$\frac{i}{h}[Q_t,\stackrel{~}{\mathrm{Op}}_{h,\stackrel{~}{h}}(a_t)]=\stackrel{~}{\mathrm{Op}}_{h,\stackrel{~}{h}}\left((H_{q_t}r_t\lambda _\lambda )a_t+\stackrel{~}{𝒪}_{\frac{1}{2}}(h\stackrel{~}{h}\lambda ^k)\right).$$
If we write $`a_t=a_t^0+𝒪_{\frac{1}{2}}(\stackrel{~}{h}\lambda ^k)`$ then
(5.47)
$$\frac{}{t}a_t^0=\left(H_{q_t}r_t(x,\xi )\lambda \frac{}{\lambda }\right)a_t^0,a_0^0amod𝒪_{\frac{1}{2}}(\stackrel{~}{h}\lambda ^k).$$
We now note that
$$(H_{q_t}r_t\lambda \frac{}{\lambda })a_t^0=H_{q_t}(a_t^0_{\lambda =(\stackrel{~}{h}/h)\xi _1}).$$
Hence, if $`K_t`$ is the transformation in $`(x,\xi ,\lambda )`$-space corresponding to $`\kappa _t`$ as in the statement of the proposition, it follows that
$$a_t^0=aK_t,$$
that the principal symbol of $`A_t`$ is $`aK_t`$ and the proposition follows. ∎
We can now define the symbol map,
$$\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^𝔭(X)A\sigma _{\mathrm{\Sigma },h,\stackrel{~}{h}}(A)S_{\mathrm{\Sigma },\frac{1}{2}}^𝔭(T^{}X)/S_{\mathrm{\Sigma },\frac{1}{2}}^𝔭^{}(T^{}X).$$
We recall from §3.1 and §3.3 that we already have the symbol map for the $`S_{\frac{1}{2}}`$ calculus:
$$\mathrm{\Psi }_{\frac{1}{2}}^{m,\stackrel{~}{m},k}(X)B\sigma _h(A)S_{\frac{1}{2}}^{m,\stackrel{~}{m},k}(T^{}X)/S_{\frac{1}{2}}^{m,\stackrel{~}{m}1,k1}(T^{}X).$$
The definition of the symbol classes (5.1) shows that
$$\begin{array}{cc}& (S_{\mathrm{\Sigma },\frac{1}{2}}^𝔭(T^{}X)𝒞^{\mathrm{}}(T^{}XU_\mathrm{\Sigma }))/S_{\mathrm{\Sigma },\frac{1}{2}}^𝔭^{}(T^{}X)=\hfill \\ & (S^{m+k_2,\stackrel{~}{m}k_2,k_1}(T^{}X)𝒞^{\mathrm{}}(T^{}XU_\mathrm{\Sigma }))/S^{m+k_2,\stackrel{~}{m}k_21,k1}(T^{}X),\hfill \end{array}$$
for any open ($`h`$-independent) neighbourhood of $`\mathrm{\Sigma }`$, $`U_\mathrm{\Sigma }`$. Hence we define
$$\sigma _{\mathrm{\Sigma },h,\stackrel{~}{h}}((I\gamma _\mathrm{\Sigma })A)\stackrel{\mathrm{def}}{=}\sigma _h((I\gamma _\mathrm{\Sigma })A),$$
and as in the proof of Lemma 5.4 we see that this definition is independent of the choice of $`\gamma _\mathrm{\Sigma }`$.
To define $`\sigma _{\mathrm{\Sigma },h,\stackrel{~}{h}}(\gamma _\mathrm{\Sigma }A)`$ we use Lemma 5.6. We choose a partition of the cut-off operator, $`\gamma _\mathrm{\Sigma }`$,
$$\underset{j=1}{\overset{J}{}}\gamma _j^2=\sigma _h(\gamma _\sigma ),\gamma _j𝒞_\mathrm{c}^{\mathrm{}}(T^{}X),$$
such that for each $`j`$, $`\mathrm{supp}\gamma _j\mathrm{\Sigma }`$ can be put into the normal form $`\mathrm{\Omega }\{\xi _1=0\}`$ by a local canonical transformation,
(5.50)
$$\begin{array}{c}\kappa _j:\mathrm{\Omega }\mathrm{\Omega }_j,(0,0)\mathrm{\Omega }T^{}^n,\\ \kappa _j(\{\xi _1=0\}\mathrm{\Omega })=\mathrm{\Sigma }\mathrm{\Omega }_j.\end{array}$$
We then choose elliptic $`h`$-Fourier Integral Operators, $`U_j`$, microlocally defined in neighbourhoods of $`\mathrm{\Omega }\times \mathrm{\Omega }_j`$ and associated to $`\kappa _j`$’s. By Proposition 5.5
$$U_j^1\mathrm{Op}_h^w(\chi _j)A\mathrm{Op}_h^w(\chi _j)U_j=\mathrm{Op}_h^w(\stackrel{~}{\chi }_j)\stackrel{~}{A}\mathrm{Op}_h^w(\stackrel{~}{\chi }_j)\mathrm{\Psi }_{\mathrm{\Sigma }^{},\frac{1}{2}}^𝔭(^n),\mathrm{\Sigma }^{}\mathrm{\Omega }=\{\xi _1=0\}.$$
In the notation of Lemma 5.6 we then define
(5.51)
$$\sigma _{\mathrm{\Sigma },h,\stackrel{~}{h}}(\gamma _\mathrm{\Sigma }A)=\underset{j=1}{\overset{J}{}}(\kappa _j^1)^{}a_{\stackrel{~}{\chi }_j}\gamma _j^2.$$
Proposition 5.7 shows that this definition is independent of the choices made here.
### 5.6. Global quantization map
From the local quantization given in §5.2 we can define a global map $`\mathrm{Op}_{\mathrm{\Sigma },h,\stackrel{~}{h}}`$. Thus let $`aS_\mathrm{\Sigma }^𝔭(T^{}X)`$ be a symbol in the class defined in (5.1).
Let $`\gamma _\mathrm{\Sigma }`$ (where we will use the same letter for the symbol and the operator) and $`V_j`$’s be as in (5.27) By shrinking $`V_2`$ if necessary we can find a finite open cover
$$V_2\underset{j=1}{\overset{J}{}}\mathrm{\Omega }_j$$
such that for each $`j`$ there exists a canonical transformation $`\kappa _j`$ satisfying (5.50).
Let $`\varphi _j`$ be a partition of unity on $`V_2`$ subordinate to the cover by $`\mathrm{\Omega }_j`$’s. Let $`a_j`$ be the unique symbol of the form
$$a_j=a_j(x,\xi _2,\mathrm{},\xi _n,\lambda ;h)$$
such that
$$\left(a_j\right)_{\lambda =(\stackrel{~}{h}/h)\xi _1}=\left(\gamma _\mathrm{\Sigma }\varphi _ja\right)\kappa _j.$$
Using the $`h`$-Fourier integral operators defined after (5.50) we put
(5.52)
$$\mathrm{Op}_{\mathrm{\Sigma },h,\stackrel{~}{h}}(a)\stackrel{\text{def}}{=}\mathrm{Op}_h^w\left((1\gamma _\mathrm{\Sigma })a\right)+\underset{j}{}U_j\stackrel{~}{\mathrm{Op}}_{h,\stackrel{~}{h}}\left(a_j\right)U_j^1.$$
In view of Proposition 5.5 we have $`\mathrm{Op}_{\mathrm{\Sigma },h,\stackrel{~}{h}}(a)\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^𝔭(X)`$.
The construction of the symbol map in §5.5 shows that
$$\sigma _{\mathrm{\Sigma },h,\stackrel{~}{h}}(\mathrm{Op}_{\mathrm{\Sigma },h,\stackrel{~}{h}}(a))a\text{mod}S_\mathrm{\Sigma }^𝔭^{}(X),𝔭^{}=(m,\stackrel{~}{m}1,k_11,k_2).$$
We recall from (5.1) that away from $`\mathrm{\Sigma }`$,
$$S_{\mathrm{\Sigma },\frac{1}{2}}^𝔭\text{ becomes }S_{\frac{1}{2}}^{m+k_2,\stackrel{~}{m}k_2,k_1},$$
and
$$S_{\mathrm{\Sigma },\frac{1}{2}}^𝔭^{}\text{ becomes }S_{\frac{1}{2}}^{m+k_2,\stackrel{~}{m}k_21,k_11}.$$
If $`\delta <1/2`$ in (5.1) then we have the usual filtration of the $`h`$-pseudodifferential calculus: near $`\mathrm{\Sigma }`$ we only gain in $`\stackrel{~}{h}`$ and away from $`\mathrm{\Sigma }`$, in $`h`$. In the case of $`\delta =1/2`$ considered here in detail we only gain in $`\stackrel{~}{h}`$, near and away $`\mathrm{\Sigma }`$.
This completes the proof of Theorem 4 and provides an explicit quantization $`\mathrm{Op}_{\mathrm{\Sigma },h,\stackrel{~}{h}}`$.
### 5.7. Approximation by finite rank operators
To estimate the number of resonances we will need to use approximation by finite rank operators.
For $`aS_{\mathrm{\Sigma },\frac{1}{2}}^{m,\stackrel{~}{m},\mathrm{},k_2}(T^{}X)`$ we need a notion of essential support. Unlike the essential support defined in §3.1 it now has to depend on $`h,\stackrel{~}{h}`$. As in , rather than introduce an equivalence class of families of sets, we will say that for an $`(h,\stackrel{~}{h})`$-dependent family of sets $`W_{h,\stackrel{~}{h}}T^{}X`$
$$\mathrm{ess}\mathrm{supp}aW_{h,\stackrel{~}{h}}$$
$$a^{}S_{\mathrm{\Sigma },\frac{1}{2}}^{m,\stackrel{~}{m},\mathrm{},k_2}(T^{}X),\mathrm{supp}a^{}W_{h,\stackrel{~}{h}},aa^{}S_{\mathrm{\Sigma },\frac{1}{2}}^{m,\mathrm{},\mathrm{},k_2}(T^{}X).$$
We notice that
$$\mathrm{ess}\mathrm{supp}aV_{h,\stackrel{~}{h}}^j,j=1,\mathrm{},N\mathrm{ess}\mathrm{supp}aV_{h,\stackrel{~}{h}}^1\mathrm{}V_{h,\stackrel{~}{h}}^N,$$
so that this formal notion of essential support behaves correctly under finite products and sums. We can now state
###### Proposition 5.10.
Suppose that $`aS_{\mathrm{\Sigma },\frac{1}{2}}^{0,0,\mathrm{},\mathrm{}}(T^{}X)`$ and
$$\mathrm{ess}\mathrm{supp}aW_{h,\stackrel{~}{h}},$$
where $`W_{h,\stackrel{~}{h}}`$ satisfies
$$W_{h,\stackrel{~}{h}}\{(x,\xi ):d((x,\xi ),\mathrm{\Sigma }W_{h,\stackrel{~}{h}})C_1h/\stackrel{~}{h}\},$$
$$W_{h,\stackrel{~}{h}}\underset{k=1}{\overset{K(h)}{}}\mathrm{exp}([1,1]H_q)(B_k),\mathrm{diam}B_kC_1(h/\stackrel{~}{h})^{\frac{1}{2}},.$$
Then for $`0<h<h_0`$, there exists a finite rank operator $`R(h)`$ such that for
$$\mathrm{Op}_{\mathrm{\Sigma },h,\stackrel{~}{h}}(a)R(h)\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^{0,\mathrm{},\mathrm{},\mathrm{}}(X),\text{rank}R(h)=C_2\stackrel{~}{h}^nK(h).$$
###### Proof.
We take an open covering of $`W_{h,\stackrel{~}{h}}`$,
$$W_{h,\stackrel{~}{h}}\underset{k=1}{\overset{K^{}(h)}{}}U_k,K^{}(h)C^{}K(h),U_k=\mathrm{exp}([1/C,1/C]H_q)V_k$$
$$\mathrm{diam}(V_k)C(h/\stackrel{~}{h})^{\frac{1}{2}},\mathrm{sup}_{(x,\xi )U_k}d((x,\xi ),U_k\mathrm{\Sigma })Ch/\stackrel{~}{h},$$
with a partition of unity on $`W_{h,\stackrel{~}{h}}`$,
$$\underset{k=1}{\overset{K^{}(h)}{}}\chi _k=1\text{ on }W_{h,\stackrel{~}{h}}\text{,}\mathrm{supp}\chi _kU_k,\chi _kS_{\mathrm{\Sigma },\frac{1}{2}}^{0,0,\mathrm{},\mathrm{}}(T^{}X).$$
If $`\psi =1_k\chi _kS_{\mathrm{\Sigma },\frac{1}{2}}^{0,0,\mathrm{},\mathrm{}}`$ then the condition on the support of $`a`$ shows that
$$\alpha ,\beta ^n,^\alpha a^\beta \psi 0.$$
Consequently the calculus of Theorem 4 gives
$$\mathrm{Op}_{\mathrm{\Sigma },h,\stackrel{~}{h}}(\psi )A\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^{0,\mathrm{},\mathrm{},\mathrm{}}(X).$$
Hence it suffices to show that for each $`k`$ there exists an operator $`R_k`$ such that
$$\mathrm{Op}_{\mathrm{\Sigma },h,\stackrel{~}{h}}(\chi _k)AR_k\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^{0,\mathrm{},\mathrm{},\mathrm{}}(X),\mathrm{rank}(R_k)C\stackrel{~}{h}^n,$$
with $`C`$ independent of $`k`$. By taking a finer partition (with a number of elements $`K^{\prime \prime }(h)C^{\prime \prime }K(h)`$) we can assume that
$$\mathrm{Op}_{\mathrm{\Sigma },h,\stackrel{~}{h}}(\chi _k)A=U_kA_kV_k$$
where $`U_k,V_k`$ are $`h`$-semiclassical Fourier Integral Operators of the form used in the construction of $`\mathrm{Op}_{\mathrm{\Sigma },h,\stackrel{~}{h}}`$, and
$$A_k=\stackrel{~}{\mathrm{Op}}_{h,\stackrel{~}{h}}(a_k)$$
$$\mathrm{supp}a_k\{(x,\xi ^{},\lambda ):|\lambda |C,|x^{}|+|\xi ^{}|C(h/\stackrel{~}{h})^{\frac{1}{2}},|x_1|<1/C\}$$
Consider commuting operators
$$Q=\left(\stackrel{~}{h}D_{x_1}\right)^2+x_1^2+(hD_{x_2})^2+x_2^2+\mathrm{}(hD_{x_n})^2+x_n^2,$$
$$Q=\stackrel{~}{\mathrm{Op}}_{h,\stackrel{~}{h}}(q),q=\lambda ^2+x_1^2+\mathrm{}\xi _n^2+x_n^2,$$
$$Q^{}=(hD_{x_2})^2+x_2^2+\mathrm{}(hD_{x_n})^2+x_n^2.$$
If $`\chi 𝒞_\mathrm{c}^{\mathrm{}}()`$, $`\chi (t)=1`$ for $`t\stackrel{~}{C}`$, $`\chi (t)=0`$ for $`t>\stackrel{~}{2}C`$, then
$$\chi (\stackrel{~}{h}Q^{}/h)\mathrm{\Psi }_{\mathrm{\Sigma }_0,\frac{1}{2}}^{0,0,0,0}(^n),$$
$$\chi (\stackrel{~}{h}Q^{}/h)\chi (Q)A_kA_k\mathrm{\Psi }_{\mathrm{\Sigma }_0,\frac{1}{2}}^{0,\mathrm{},\mathrm{},\mathrm{}}.$$
The standard analysis of the spectrum of harmonic oscillators shows that $`\chi (\stackrel{~}{h}Q^{}/h)\chi (Q)`$ is a finite rank operator and its rank is bounded by $`C^{}\stackrel{~}{h}^n`$. Hence we can take $`R_k=\chi (Q)\chi (\stackrel{~}{h}Q^{}/h)A_k`$. ∎
## 6. General upper bounds in regions of size $`h`$
### 6.1. Bound for the number of eigenvalues of a self-adjoint operator
With the calculus developed in §5 we can follow the standard procedure of modifying the operator near the energy surface, now at the limiting scale. For that we introduce the second small parameter $`\stackrel{~}{h}`$ which eventually will be fixed, as $`h0`$.
Let $`\chi 𝒞_\mathrm{c}^{\mathrm{}}(;[0,1])`$ be equal to one near $`0`$. We then define
(6.1)
$$a(x,\xi ,h)\stackrel{\mathrm{def}}{=}\chi \left(\frac{\stackrel{~}{h}p(x,\xi )}{h}\right).$$
Then, in terms of Definition 5.1, $`aS_{\mathrm{\Sigma },0}^{0,0,\mathrm{},\mathrm{}}(T^{}X)`$, $`\mathrm{\Sigma }=p^1(0)`$. Although this class of symbols corresponds to a class $`\mathrm{\Psi }_{\mathrm{\Sigma },0}^{}`$ we will use the larger class $`\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^{}`$ since it was presented in detail in §5. We stress that this is done for convenience only and an examination of §5 shows how the simpler calculus is constructed without the $`S_{\frac{1}{2}}`$ complications.
The operator
$$\stackrel{~}{P}\stackrel{\mathrm{def}}{=}P+i(h/\stackrel{~}{h})Az,A\stackrel{\mathrm{def}}{=}\stackrel{~}{\mathrm{Op}}_{h,\stackrel{~}{h}}(a),$$
is elliptic in $`\mathrm{\Psi }_\mathrm{\Sigma }^{0,0,2,1}(X)`$, in the sense that
$$|\sigma _{\mathrm{\Sigma },h,\stackrel{~}{h}}(\stackrel{~}{P})|>C\left(d(,\mathrm{\Sigma })+h/\stackrel{~}{h}\right).$$
This is most clearly seen locally when $`p=\xi _1`$ and $`\lambda =(h/\stackrel{~}{h})^1\xi _1`$:
$$\stackrel{~}{p}=\frac{h}{\stackrel{~}{h}}\left(\lambda +i\chi \left(\lambda \right)\right),|\stackrel{~}{p}|>C\frac{h}{\stackrel{~}{h}}\lambda .$$
Using Theorem 4 we can construct a parametrix for $`\stackrel{~}{P}z`$, $`|z|Ch`$, uniformly in $`h`$, which for $`\stackrel{~}{h}`$ small enough (keeping $`\stackrel{~}{h}M`$ large and constant) gives an exact inverse:
$$(\stackrel{~}{P}z)^1=𝒪\left(\frac{\stackrel{~}{h}}{h}\right):L^2(X)L^2(X),|z|Ch.$$
Proposition 5.10 gives a finite rank operator $`R`$ such that
$$A=R+E,\mathrm{rank}(R)C\stackrel{~}{h}^nh^{n+1},E\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^{0,\mathrm{},\mathrm{},\mathrm{}},$$
so that in particular, by Proposition 5.5,
$$E_{L^2L^2}=𝒪(\stackrel{~}{h}^{\mathrm{}}).$$
Now we write
$$\begin{array}{cc}\hfill Pz& =(\stackrel{~}{P}z)(Ii(h/\stackrel{~}{h})(\stackrel{~}{P}z)^1(R+E))\hfill \\ & =(\stackrel{~}{P}z)(Ii(h/\stackrel{~}{h})(\stackrel{~}{P}z)^1E)(I+K(z)),\hfill \end{array}$$
where
$$K(z)=i(h/\stackrel{~}{h})(I+i(h/\stackrel{~}{h})(\stackrel{~}{P}z)^1E)^1(\stackrel{~}{P}z)^1R,$$
$$\mathrm{rank}(K(z))Mh^{n+1},K(z)=𝒪(1):L^2(X)L^2(X).$$
This implies that
$$h(z)\stackrel{\mathrm{def}}{=}det(I+K(z))=𝒪(\mathrm{exp}(CMh^{n+1})),|z|C_0h,$$
and that the zeros of $`h(z)`$ are the eigenvalues of $`P`$ in $`|z|Ch`$.
The bound on the number of eigenvalues will follow standard estimates The estimate we need is this: if $`\mathrm{log}|h(z)|K`$ in $`R_1`$, where $`R_1`$ is a rectangle, and $`R_2`$ is another rectangle strictly inside $`R_2`$, $`\mathrm{log}|h(z_0)|>K`$, for some $`z_0R_2`$, then the number of zeros of $`h(z)`$ in $`R_2`$ is bounded by $`C(R_1,R_2)K`$, with a dilation invariant constant, $`C(tR_1,tR_2)=C(R_1,R_2)`$. once we show that
$`|h(z_0)|>\mathrm{exp}(CMh^{n+1})`$ at some $`z_0`$, $`|z_0|C_1h`$, $`C_1<C_0`$.
For that we take $`z_0`$ with $`|Imz_0|>C_1h`$ so that, by self-adjointness,
$$(Pz_0)^1=𝒪((C_1h)^1):L^2(X)L^2(X).$$
We then see that
$$(I+K(z_0))^1=I+L(z_0),$$
where
$$I+L(z)=(I+i(h/\stackrel{~}{h})(Pz)^1(R+E))(I+i(h/\stackrel{~}{h})(\stackrel{~}{P}z)^1E),|Imz|>C_1h.$$
The operator $`L(z_0)`$ is of trace class and using the rank of $`R`$ we have the estimate
$$h(z_0)^1=det(I+L(z_0))=𝒪(\mathrm{exp}(CMh^{n+1})),$$
which shows that
(6.2)
$$|\mathrm{Spec}PD(0,h)|=𝒪(h^{n+1}).$$
### 6.2. Proof of Theorem 2
To establish the estimate (1.16) we proceed as in the case of eigenvalues but using the scaled operator $`P_\theta `$ instead – see §3.5. In this section we take
$$\theta =C_0h$$
for $`C_0`$, a large and fixed constant. We recall from §3.1 that
$$X=X_0(^nB(0,R_0)),$$
We can also have more neighbourhoods of infinity (see §3.1) but for simplicity of notation we restrict ourselves to the case above.
The operator $`P_\theta `$ can be written as $`P_\theta =P_1+iP_2`$, where $`P_j`$ are formally self-adjoint on $`L^2(X)`$. We consider the Weyl symbols, $`p_1`$ and $`p_2`$ (defined, we recall, modulo $`𝒪(h^2)`$) of these two operators. In view of (3.30) we have
(6.5)
$$\begin{array}{c}p_1_{T_{X_0^nB(0,2R_0)}^{}X}=p_{T_{X_0^nB(0,2R_0)}^{}X},|p_1p|c_1h\xi ^2,\\ p_2_{X_0}=0,3C_1h|\xi |^2p_2_{T_{^nB(0,2R_0)}^{}X}C_1h|\xi |^2,|p_2|c_1h\xi ^2.\end{array}$$
We also note that our assumptions give
$$|p_1|\delta \xi |\xi |,$$
with some small fixed $`\delta `$. Now let $`\chi 𝒞_\mathrm{c}^{\mathrm{}}(,[0,1])`$, satisfy
$$\chi (t)=\{\begin{array}{cc}1\hfill & |t|1\hfill \\ 0\hfill & |t|2\hfill \end{array}$$
With this $`\chi `$ we define
$$a(x,\xi ,h)\stackrel{\mathrm{def}}{=}\chi \left(\frac{\stackrel{~}{h}p(x,\xi )}{h}\right)\chi \left(\frac{|x|}{3R_0}\right),$$
where $`\stackrel{~}{h}`$ is small,
$$\stackrel{~}{h}^1C_1c_1C_0,$$
in (6.5) but eventually fixed.
We now choose a compact $`𝒞^{\mathrm{}}`$ hypersurface $`\mathrm{\Sigma }`$ so that
$$\mathrm{\Sigma }T_{X_0B(0,6R_0)B(0,R_0)}^{}X=\{(x,\xi ):p(x,\xi )=0\}.$$
It follows that
$$a(x,\xi ,h)S_\mathrm{\Sigma }^{0,0,\mathrm{},0}(T^{}X),$$
and continuing in the same spirit as in §6.1, we put
$$A\stackrel{\mathrm{def}}{=}\stackrel{~}{\mathrm{Op}}_{h,\stackrel{~}{h}}(a).$$
On the level of symbols we have
(6.6)
$$|p_1|(h/\stackrel{~}{h})p_2(x,\xi )+(h/\stackrel{~}{h})a(x,\xi ,h)(h/\stackrel{~}{h})/C.$$
Hence if we put
$$\stackrel{~}{P}=P_\theta i(h/\stackrel{~}{h})A,$$
then for $`|Rez|Ch`$, $`ImzCh`$,
$$(\stackrel{~}{h}/h)(\stackrel{~}{P}z)\mathrm{\Psi }_\mathrm{\Sigma }^{0,0,2,1}(X).$$
We claim that $`\stackrel{~}{P}z`$ is invertible for $`zD(0,Ch)`$ and
(6.7)
$$(\stackrel{~}{P}z)^1=𝒪(\stackrel{~}{h}/h).$$
In fact, let $`WT^{}X`$ be set in which $`(\stackrel{~}{h}/h)(\stackrel{~}{p}z)`$ is elliptic in $`S_\mathrm{\Sigma }^{0,0,2,1}(T^{}X)`$. The estimate (6.6) shows that $`Im(\stackrel{~}{p}z)C(h/\stackrel{~}{h})`$ on a neighbourhood of $`\mathrm{}W\{|Re(\stackrel{~}{p}z)|<\delta \}`$. Now, let $`\mathrm{\Psi }_1`$ and $`\mathrm{\Psi }_2`$ be as in Lemma 3.2, with $`\mathrm{ess}\mathrm{supp}\mathrm{\Psi }_1W`$. Then proceeding as in §4.1 we obtain
$$(\stackrel{~}{P}z)\mathrm{\Psi }_2u(h/\stackrel{~}{h})\mathrm{\Psi }_2u/C𝒪(h^{\mathrm{}})u,u𝒞_\mathrm{c}^{\mathrm{}}(X).$$
Ellipticity of $`(\stackrel{~}{h}/h)(\stackrel{~}{P}z)`$ in $`W`$ shows that
$$(\stackrel{~}{h}/h)(\stackrel{~}{P}z)\mathrm{\Psi }_1u\mathrm{\Psi }_1u/C𝒪(\stackrel{~}{h}^{\mathrm{}}),$$
that is
$$(\stackrel{~}{P}z)\mathrm{\Psi }_1u(h/\stackrel{~}{h})\mathrm{\Psi }_1u/C𝒪(h\stackrel{~}{h}^{\mathrm{}}).$$
Lemma 3.3 then gives (6.7).
We can now proceed as in §6.1 and obtain the bound on the number of resonances (1.16). The complex analytic argument outlined in the footnote is used in a rectangle
$$[Ch,Ch]+i[C_2h,Ch]$$
where $`C_2`$ is large enough to guarantee a lower bound for $`Im(P_\theta z)`$ when $`ImzC_2h`$. We can still take $`C_2`$ proportional to the other large constants.
## 7. The escape function for hyperbolic flows and its $`h`$ dependent regularizations
In this section we modify \[29, Sect.5\] and construct a regularized escape function depending on a small parameter, essentially $`h/\stackrel{~}{h}`$. We recall that we assume that $`p𝒞^{\mathrm{}}(T^{}X;)`$ satisfies
(7.3)
$$\begin{array}{c}p=0dp0\\ |x|R,|p(x,\xi )|<2\delta \mathrm{exp}tH_p(x,\xi )\mathrm{}\text{ for either }t\mathrm{}\text{ or }t\mathrm{}\text{.}\end{array}$$
We also recall the result of \[11, Appendix\]:
###### Proposition 7.1.
Suppose that (7.3) holds and that $`\widehat{K}`$ is the trapped set,
(7.4)
$$\widehat{K}\stackrel{\mathrm{def}}{=}\{\rho T^{}X:\mathrm{exp}(tH_p)(\rho )\to ̸\mathrm{},t\pm \mathrm{},|p(\rho )|\delta \}T^{}X.$$
Then for any two neighbourhoods, $`U,V`$, of $`\widehat{K}`$, $`\overline{U}V`$ there exists $`G_0𝒞^{\mathrm{}}(T^{}X)`$ such that
(7.7)
$$\begin{array}{c}\mathrm{supp}G_0T^{}XU,H_pG_00,H_pG_{p^1([2\delta ,2\delta ])}C,\\ H_pG_0_{p^1([\delta ,\delta ])V}1.\end{array}$$
### 7.1. Dynamical assumptions
We start with the hyperbolicity assumptions \[29, §5\] weaker than the more standard assumptions in §1. Let $`\widehat{K}`$ be the compact trapped set near zero energy given by (7.4). The trapped set at zero energy is given by $`K=\widehat{K}p^1(0)`$. We also have $`\widehat{K}=\widehat{\mathrm{\Gamma }}_+\widehat{\mathrm{\Gamma }}_{}`$, where
(7.8)
$$\widehat{\mathrm{\Gamma }}_\pm \stackrel{\mathrm{def}}{=}\{(x,\xi )T^{}X:|p(x,\xi )|\delta ,\mathrm{exp}(tH_p)(x,\xi )\to ̸\mathrm{},t\mathrm{}\},$$
and the sets $`\widehat{K}`$, $`\widehat{\mathrm{\Gamma }}_\pm `$ are clearly invariant under the flow,
(7.9)
$$\mathrm{exp}(tH_p)(\widehat{K})\widehat{K},\mathrm{exp}(tH_p)(\widehat{\mathrm{\Gamma }}_\pm )\widehat{\mathrm{\Gamma }}_\pm .$$
We can now state the dynamical hypothesis.
* In a neighbourhood, $`\mathrm{\Omega }_{\rho _0}`$ of any $`\rho _0K`$,
$$\widehat{\mathrm{\Gamma }}_\pm =\underset{\rho \mathrm{\Omega }_{\rho _0}\widehat{\mathrm{\Gamma }}_\pm }{}\widehat{\mathrm{\Gamma }}_{\pm ,\rho },\rho \widehat{\mathrm{\Gamma }}_{\pm ,\rho },$$
$$\widehat{\mathrm{\Gamma }}_{\pm ,\rho }\widehat{\mathrm{\Gamma }}_{\pm ,\rho ^{}}=\mathrm{},\text{ or }\widehat{\mathrm{\Gamma }}_{\pm ,\rho }=\widehat{\mathrm{\Gamma }}_{\pm ,\rho ^{}}.$$
* Each $`\widehat{\mathrm{\Gamma }}_{\pm ,\rho }`$ is a closed $`𝒞^1`$ manifold of dimension $`n+d`$, with $`d0`$ fixed, and the dependence
$$\mathrm{\Omega }_{\rho _0}\widehat{\mathrm{\Gamma }}_\pm \rho T_\rho \widehat{\mathrm{\Gamma }}_{\pm ,\rho }$$
is continuous.
* If $`E_\rho ^\pm \stackrel{\mathrm{def}}{=}T_\rho \widehat{\mathrm{\Gamma }}_{\pm ,\rho }`$, then $`E_\rho ^++E_\rho ^{}=T_\rho p^1(p(\rho ))T_\rho (T^{}X)`$, $`H_p(\rho )E_\rho ^\pm `$, and
(7.10)
$$d(\mathrm{exp}tH_p)_\rho (X)Ce^{\pm \lambda t}X,\rho K,\text{ for all }XT_\rho (T^{}X)/E_\rho ^{}\text{}t0\text{.}$$
The above definition makes sense since by (7.9) $`d(\mathrm{exp}tH_p)_\rho (E_\rho ^\pm )=E_{\mathrm{exp}tH_p(\rho )}`$, $`\rho \widehat{\mathrm{\Gamma }}_\pm `$, we have
$$d(\mathrm{exp}tH_p)_\rho T_\rho (T^{}X)/E_\rho ^{}T_{\mathrm{exp}tH_p(\rho )}(T^{}X)/E_{\mathrm{exp}tH_p(\rho )}^{},\rho K,$$
and we choose continuously dependent norms in the last estimate in (7.10). We also note that $`XT_\rho (T^{}X)/E_\rho ^{}`$ implies that $`X`$ can be identified with a vector tangent to $`p^1(p(\rho ))`$.
In \[29, §5\] it is shown that there exist two functions, $`\phi _\pm 𝒞^{1,1}(T^{}X)`$, $`\phi _\pm 0`$, $`H_p^k\phi _\pm 𝒞^{1,1}(T^{}X),k`$, such that for $`\rho `$ in a small neighbourhood of $`K`$,
$$\begin{array}{cc}\hfill H_p\phi _\pm (\rho )& \phi _\pm (\rho ),H_p^k\phi _\pm (\rho )=𝒪(\phi _\pm (\rho )),k,,\hfill \\ \hfill \phi _\pm (\rho )& d(\rho ,\widehat{\mathrm{\Gamma }}_\pm ),\phi _+(\rho )+\phi _{}(\rho )d(\rho ,\widehat{K})^2,\hfill \end{array}$$
and where $`d(,\mathrm{\Gamma })`$ is the distance to a closed set $`\mathrm{\Gamma }`$. The notation $`fg`$, means that there exists a constant $`C>0`$ such
$$0g/CfCg.$$
The simple model (2.3) is given in §2. Here we modify the construction to obtain suitably regularized functions $`\widehat{\phi }_\pm `$
### 7.2. Regularization of $`\phi _\pm `$.
We start with two general lemmas:
###### Lemma 7.2.
Suppose $`\mathrm{\Gamma }^m`$ is a closed set. For any $`ϵ>0`$ there exists $`\phi _ϵ𝒞^{\mathrm{}}(^m)`$ such that
$$\phi _ϵϵ,\phi _ϵd(,\mathrm{\Gamma })^2+ϵ,^\alpha \phi _ϵ=𝒪(\phi _ϵ^{1|\alpha |/2}),$$
uniformly on compact sets.
###### Proof.
We can find a sequence $`x_j^m`$ such that
$$\underset{j}{}B(x_j,d(x_j,\mathrm{\Gamma })/8)=^m\mathrm{\Gamma },$$
every $`xQ\mathrm{\Gamma }`$, $`Q^m`$, is in at most $`N_0=N_0(Q)`$ balls $`B(x_j,d(x_j,\mathrm{\Gamma })/2)`$.
Let $`\chi 𝒞_\mathrm{c}^{\mathrm{}}(^m;[0,1])`$ be supported in $`B(0,1/4)`$, and be identically one in $`B(0,1/8)`$. We define
$$\phi _ϵ(x)\stackrel{\mathrm{def}}{=}ϵ+\underset{d(x_j,\mathrm{\Gamma })>\sqrt{ϵ}}{}d(x_j,\mathrm{\Gamma })^2\chi \left(\frac{xx_j}{d(x_j,\mathrm{\Gamma })+\sqrt{ϵ}}\right)$$
We first note that the number non-zero terms in the sum is uniformly bounded by $`N_0`$. In fact, $`d(x_j,\mathrm{\Gamma })+\sqrt{ϵ}<2d(x_j,\mathrm{\Gamma })`$, and hence if $`\chi ((xx_j)/(d(x_j,\mathrm{\Gamma })+\sqrt{ϵ}))0`$ then
$$1/4|xx_j|/(d(x_j,\mathrm{\Gamma })+\sqrt{ϵ})(1/2)|xx_j|/d(x_j,\mathrm{\Gamma }),$$
and $`xB(x_j,d(x_j,d(x_j,\mathrm{\Gamma }))/2)`$. This shows that $`\phi _ϵ(x)2N_0(ϵ+d(x,\mathrm{\Gamma })^2)`$, and
$$^\alpha \phi _ϵ(x)=𝒪((d(x,\mathrm{\Gamma })^2+ϵ)^{1|\alpha |/2}),$$
uniformly on compact sets.
To see the lower bound on $`\phi _ϵ`$ we first consider the case when $`d(x,\mathrm{\Gamma })C\sqrt{ϵ}`$.
$$\phi _ϵ(x)ϵ(ϵ+d(x,\mathrm{\Gamma })^2)/C^{}.$$
If $`d(x,\mathrm{\Gamma })>C\sqrt{ϵ}`$ then for at least one $`j`$, $`\chi ((xx_j)/(d(x_j,\mathrm{\Gamma })+\sqrt{ϵ}))=1`$ (since the balls $`B(x_j,d(x_j,\mathrm{\Gamma })/8)`$ cover the complement of $`\mathrm{\Gamma }`$, and $`\chi (t)=1`$ if $`|t|1/8`$). Thus
$$\phi _ϵ(x)ϵ+d(x_j,\mathrm{\Gamma })^2(ϵ+d(x,\mathrm{\Gamma })^2)/C,$$
which concludes the proof. ∎
For future use we also record the following
###### Lemma 7.3.
Suppose $`\phi 𝒞^{1,1}(^m)`$, $`\phi 0`$, and for a vectorfield $`V𝒞^{\mathrm{}}(^m;^m)`$, $`V^k\phi =𝒪(\phi )`$, $`V^k\varphi 𝒞^{1,1}(^m)`$, $`k`$. Then, uniformly on compact sets,
$$dV^k\phi =𝒪(\phi ^{\frac{1}{2}}),k.$$
###### Proof.
For some $`C>0`$ the $`𝒞^{1,1}`$ function $`C\phi V^k\phi `$ is non-negative. Hence using the standard estimate based on Taylor’s formula,
$$|d\phi |^2=𝒪(\phi ),|d(C\phi V^k\phi )|^2=𝒪(C\phi V^k\phi )=𝒪(\phi ).$$
The lemma follows. ∎
We now have
###### Proposition 7.4.
Let $`\widehat{\mathrm{\Gamma }}_\pm `$ be given by (7.8). For any small $`ϵ>0`$ there exist functions $`\widehat{\phi }_\pm 𝒞^{\mathrm{}}(T^{}X;[0,\mathrm{}))`$ such that in a neighbourhood of $`\widehat{K}`$,
(7.11)
$$\begin{array}{cc}\hfill \widehat{\phi }_\pm (\rho )& d(\rho ,\widehat{\mathrm{\Gamma }}_\pm )^2+Cϵ,\hfill \\ \hfill H_p\widehat{\phi }_\pm (\rho )+Cϵ& \widehat{\phi }_\pm (\rho ),\hfill \\ \hfill ^\alpha H_p^k\widehat{\phi }_\pm (\rho )& =𝒪(\widehat{\phi }_\pm (\rho )^{1|\alpha |/2}),k,\hfill \\ \hfill \widehat{\phi }_+(\rho )+\widehat{\phi }_{}(\rho )& d(\rho ,\widehat{K})^2+Cϵ.\hfill \end{array}$$
###### Proof.
We modify the arguments of \[29, §5\], roughly speaking, adding an $`𝒪(ϵ)`$ error to all the estimates. Let $`\phi _\pm `$ be the functions obtained using Lemma 7.2 with $`\mathrm{\Gamma }=\mathrm{\Gamma }_\pm `$. We now put
$$\widehat{\phi }_\pm (\rho )\stackrel{\mathrm{def}}{=}_{}g_T(t)\phi _\pm (\mathrm{exp}tH_p(\rho ))𝑑t,$$
$$g_T𝒞_\mathrm{c}^{\mathrm{}}((1,T+1)),\mathrm{supp}g_T^{}[1,1][T1,T+1],$$
$$g_T^{}_{[1,1]}0,g_T^{}_{[T1,T+1]}0,g_T^{}(0)=1,g_T^{}(T)=1.$$
To check (7.11) we note that, by definition, $`\phi _\pm (\rho )d(\rho ,\widehat{\mathrm{\Gamma }}_\pm )^2+Cϵ`$. The assumptions (7.10) imply (see \[29, Lemma 5.2\]) that
$$C,T0,\mathrm{\Omega }_TK,\text{an open set, }d(\mathrm{exp}(\pm TH_p)(\rho ),\widehat{\mathrm{\Gamma }}_\pm )Ce^{T/C}d(\rho ,\widehat{\mathrm{\Gamma }}_\pm ).$$
Hence, with constants depending on $`T`$,
$$\widehat{\phi }_+(\rho )\phi _+(\mathrm{exp}(TH_p)(\rho ))\phi _+(\rho )d(\rho ,\mathrm{\Gamma }_+)^2+Cϵ,$$
$$\widehat{\phi }_{}(\rho )\phi _{}(\rho )d(\rho ,\mathrm{\Gamma }_{})^2+Cϵ.$$
This shows the first statement in (7.11).
The assumptions on $`g_T`$ also show that
$$H_p\widehat{\phi }_\pm (\rho )\phi _\pm (\mathrm{exp}TH_p(\rho ))\phi _\pm (\rho )d(\mathrm{exp}TH_p(\rho ),\widehat{\mathrm{\Gamma }}_\pm )^2d(\rho ,\widehat{\mathrm{\Gamma }}_\pm )^2+𝒪(ϵ).$$
so that for $`T`$ large enough and for $`\rho `$ in a small neighbourhood of $`K`$, (again with $`T`$ depenendent constants)
$$H_p\widehat{\phi }_\pm (\rho )+Cϵd(\rho ,\widehat{\mathrm{\Gamma }}_\pm )^2+C^{}ϵ\widehat{\phi }_\pm (\rho ).$$
This proves the second part of (7.11). The third part is proved using Lemma 7.3 for $`|\alpha |=1`$ and the estimates on $`\phi _\pm `$ in general.
To prove the last statement in (7.11) we first see that the transversality, $`E_{\rho _0}^++E_{\rho _0}^{}=T_{\rho _0}(T^{}X)`$, and the continuity, $`\rho E_\rho ^\pm `$, assumed in (7.10) imply that for $`\rho ,\rho _1,\rho _2,`$ near a point $`\rho _0K`$,
$$d(\rho ,\widehat{\mathrm{\Gamma }}_{+,\rho _1}\widehat{\mathrm{\Gamma }}_{,\rho _2})d(\rho ,\widehat{\mathrm{\Gamma }}_{+,\rho _1})+d(\rho ,\widehat{\mathrm{\Gamma }}_{,\rho _2}).$$
Hence
$$\begin{array}{cc}\hfill \widehat{\phi }_+(\rho )+\widehat{\phi }_{}(\rho )+𝒪(ϵ)& d(\rho ,\widehat{\mathrm{\Gamma }}_+)^2+d(\rho ,\widehat{\mathrm{\Gamma }}_{})^2+Cϵ\hfill \\ & d(\rho ,\widehat{\mathrm{\Gamma }}_{+,\rho ^{}})^2+d(\rho ,\widehat{\mathrm{\Gamma }}_{,\rho ^{}})^2+Cϵ\hfill \\ & d(\rho ,\widehat{\mathrm{\Gamma }}_{+,\rho ^{}}\widehat{\mathrm{\Gamma }}_{,\rho ^{}})^2+Cϵ.\hfill \end{array}$$
If we choose $`\rho ^{}K`$ so that $`d(\rho ,\widehat{K})=d(\rho ,\rho ^{})`$ then
$$d(\rho ,\widehat{\mathrm{\Gamma }}_{+,\rho ^{}}\widehat{\mathrm{\Gamma }}_{,\rho ^{}})^2d(\rho ,\rho ^{})^2=d(\rho ,\widehat{K})^2,$$
proving that
$$\widehat{\phi }_+(\rho )+\widehat{\phi }_{}(\rho )d(\rho ,\widehat{K})^2+𝒪(ϵ).$$
The opposite inequality is obtained by choosing $`\rho _\pm \widehat{\mathrm{\Gamma }}_\pm `$ such that $`d(\rho ,\rho _\pm )=d(\rho ,\widehat{\mathrm{\Gamma }}_\pm )`$. Then using the transversality of $`\widehat{\mathrm{\Gamma }}_+`$, $`\widehat{\mathrm{\Gamma }}_{}`$
$$\begin{array}{cc}\hfill d(\rho ,\widehat{K})^2& d(\rho ,\widehat{\mathrm{\Gamma }}_{+,\rho _+}\widehat{\mathrm{\Gamma }}_{,\rho _{}})^2d(\rho ,\widehat{\mathrm{\Gamma }}_{+,\rho _+})^2+d(\rho ,\widehat{\mathrm{\Gamma }}_{,\rho _{}})^2\hfill \\ & d(\rho ,\rho _+)^2+d(\rho ,\rho _{})^2=d(\rho ,\widehat{\mathrm{\Gamma }}_+)^2+d(\rho ,\widehat{\mathrm{\Gamma }}_{})^2\hfill \\ & \widehat{\phi }_+(\rho )+\widehat{\phi }_{}(\rho )+𝒪(ϵ).\hfill \end{array}$$
### 7.3. Regularized escape function
We now use the functions constructed in Proposition 7.4 to obtain an escape function near $`K`$. We first need the following
###### Lemma 7.5.
Then for $`|\alpha |+k1`$ we have
$$_\rho ^\alpha H_p^k\mathrm{log}(\widehat{\phi }_\pm )=𝒪(\widehat{\phi }_\pm ^{\frac{|\alpha |}{2}})$$
###### Proof.
Let $`f(t)=\mathrm{log}(t)`$. Then
$$f^{(k)}(\widehat{\phi }_\pm )=𝒪\left(\frac{1}{\widehat{\phi }_\pm ^k}\right),k1,$$
and for $`|\alpha |+k1`$, $`_\rho ^\alpha H_p^kf(\widehat{\phi }_\pm )`$ is a finite linear combination of terms
$$f^{(l)}(\widehat{\phi }_\pm )\left(_\rho ^{\alpha _1}H_p^{k_1}\widehat{\phi }_\pm \right)\mathrm{}\left(_\rho ^\alpha _{\mathrm{}}H_p^k_{\mathrm{}}\widehat{\phi }_\pm \right)=𝒪(1)\underset{j=1}{\overset{\mathrm{}}{}}\frac{_\rho ^{\alpha _j}H_p^{k_j}\widehat{\phi }_\pm }{\widehat{\phi }_\pm },$$
with
$$|\alpha _j|+k_j1,\alpha _1+\mathrm{}+\alpha _{\mathrm{}}=\alpha ,k_1+\mathrm{}+k_{\mathrm{}}=k.$$
The estimates in (7.11) show that $`_\rho ^{\alpha _j}H_p^{k_j}\widehat{\phi }_\pm /\widehat{\phi }_\pm =𝒪(\widehat{\phi }_\pm ^{|\alpha _j|/2})`$, and hence
$$_\rho ^\alpha H_p^kf(\widehat{\phi }_\pm )=𝒪(\widehat{\phi }_\pm ^{\frac{|\alpha |}{2}}),$$
proving the lemma. ∎
We are now ready for the main results of this section.
###### Lemma 7.6.
Let $`\widehat{\phi }_\pm `$ be given in Proposition 7.4 and
(7.12)
$$\widehat{G}\stackrel{\mathrm{def}}{=}\left(\mathrm{log}(Mϵ+\widehat{\phi }_{})\mathrm{log}(Mϵ+\widehat{\phi }_+)\right).$$
Then in a neighbourhood of $`K`$ we have
(7.15)
$$\begin{array}{c}_\rho ^\alpha H_p^k\widehat{G}=𝒪_M(\mathrm{min}(\widehat{\phi }_+,\widehat{\phi }_{})^{\frac{|\alpha |}{2}})=𝒪_M(ϵ^{\frac{|\alpha |}{2}}),|\alpha |+k1,\\ d(\rho ,\widehat{K})^2CϵH_p\widehat{G}1/C,\end{array}$$
where, for the second estimate, $`M`$ has to be chosen large enough, independently of $`ϵ`$, and $`C`$ is a large constant.
###### Proof.
We observe that, with constants depending on $`M`$, $`\widehat{\phi }_\pm +Mϵ`$ has the same properties as $`\widehat{\phi }_\pm `$. Hence the estimates on $`_\rho ^\alpha H_p^k\widehat{G}`$ follow directly from the definition (7.12) and from Lemma 7.5. To check the second part of (7.15) we compute, using Proposition 7.4,
$$H_p\widehat{G}=\left(\frac{H_p\widehat{\phi }_{}}{\widehat{\phi }_{}+Mϵ}\frac{H_p\widehat{\phi }_+}{\widehat{\phi }_++Mϵ}\right)\frac{1}{C_1}\left(\frac{\widehat{\phi }_{}C_2ϵ}{\widehat{\phi }_{}+Mϵ}+\frac{\widehat{\phi }_+C_2ϵ}{\widehat{\phi }_++Mϵ}\right).$$
From (7.11) we also have
$$d(\rho ,\widehat{K})^2Cϵ\mathrm{max}(\widehat{\phi }_+,\widehat{\phi }_{})(C/2𝒪(1))ϵ>C_3ϵ,$$
where $`C_3`$ can be as large as we like depending on the choice of $`C`$. Hence, since $`x(xC_2)/(x+M)`$ is increasing,
$$H_p\widehat{G}\frac{1}{C_1}\left(\frac{C_3C_2}{C_3+M}\frac{C_2}{M}\right)\frac{1}{C},$$
if we choose $`C_3MC_2`$. ∎
We now modify $`\widehat{G}`$ using $`G_0`$ given in Proposition 7.1:
###### Proposition 7.7.
Let us fix $`\delta _0>0`$. Then there exist $`\widehat{\chi },\chi _0𝒞_\mathrm{c}^{\mathrm{}}(T^{}X)`$, $`C_0>0`$, and a neighbourhoood $`V`$ of $`K`$, such that
$$G\stackrel{\mathrm{def}}{=}\widehat{\chi }\widehat{G}+C_0\left(\mathrm{log}\frac{1}{ϵ}\right)\chi _0G_0,$$
satisfies
(7.22)
$$\begin{array}{c}^\alpha H_p^kG=\{\begin{array}{cc}𝒪(\mathrm{log}(1/ϵ))\hfill & \alpha =0\hfill \\ 𝒪(ϵ^{|\alpha |/2})\hfill & \text{otherwise}\hfill \end{array},\\ d(\rho ,\widehat{K})^2Cϵ,\rho VH_pG(\rho )1/C,\\ \rho p^1([\delta ,\delta ])V,|x(\rho )|3R_0H_pG(\rho )\mathrm{log}(1/ϵ),\\ H_pG(\rho )\delta _0\mathrm{log}(1/ϵ),\rho T^{}X.\end{array}$$
In addition we have
(7.23)
$$\frac{\mathrm{exp}G(\rho )}{\mathrm{exp}G(\mu )}C_0\frac{\rho \mu }{\sqrt{ϵ}}^{N_0},$$
for some constants $`C_0`$ and $`N_0`$.
###### Proof.
We obtain $`G_0`$ from Proposition 7.1 taking for $`V`$ a neighbourhood of $`\widehat{K}`$ in which the estimates of Lemma 7.6 hold. We have $`^\alpha H_p^kG_0=𝒪_{k,|\alpha |}(1)`$, and consequently for any $`\chi _0𝒞_\mathrm{c}^{\mathrm{}}(T^{}X)`$,
$$^\alpha H_p^k\left(\mathrm{log}(1/ϵ)\chi _0G_0\right)=𝒪_{k,|\alpha |}(\mathrm{log}(1/ϵ))=\{\begin{array}{cc}𝒪(\mathrm{log}(1/ϵ))\hfill & \alpha =0\hfill \\ 𝒪(ϵ^{|\alpha |/2})\hfill & \text{otherwise}\hfill \end{array}.$$
From Lemma 7.6 we obtain, again for any $`\widehat{\chi }𝒞_\mathrm{c}^{\mathrm{}}(T^{}X)`$,
$$^\alpha H_p^k(\widehat{\chi }\widehat{G})=\{\begin{array}{cc}𝒪(\mathrm{log}(1/ϵ))\hfill & \alpha =0\hfill \\ 𝒪(ϵ^{|\alpha |/2})\hfill & \text{otherwise}\hfill \end{array}.$$
The loss compared to (7.15) is due to the presence of the cut-off function.
We take $`\chi _0𝒞_\mathrm{c}^{\mathrm{}}(T^{}X;[0,1])`$ to be identically equal to $`1`$ in
$$p^1([\delta ,\delta ])\{(x,\xi ):|x|3R_0\}.$$
For $`\widehat{\chi }𝒞_\mathrm{c}^{\mathrm{}}(T^{}X)`$ we take a function which is supported in a neighbhourhood of $`\widehat{K}`$ where (7.15) holds, and identically $`1`$ in $`V`$. Hence for $`\rho p^1([\delta ,\delta ])V`$, $`|x(\rho )|3R_0`$,
$$H_pG(\rho )=C_0\mathrm{log}(1/ϵ)H_pG_0(\rho )+H_p(\widehat{\chi }\widehat{G})(\rho )C_0\mathrm{log}(1/ϵ)𝒪(1)\mathrm{log}(1/ϵ)\mathrm{log}(1/ϵ),$$
if $`C_0`$ is taken large enough. For $`\rho V`$, $`\widehat{\chi }(\rho )=1`$, and
$$H_pG(\rho )=C_0\mathrm{log}(1/ϵ)H_pG_0(\rho )+H_p\widehat{G}(\rho )H_p\widehat{G}(\rho ),$$
and if $`d(\rho ,\widehat{K})Cϵ`$, $`H_pG(\rho )1/C`$. To complete the proof of (7.22) we need to define $`\chi _0`$ for $`|x|R_0`$. That is essentially done as in §4.2 where it was based on Lemma 4.1. Let $`T`$ and $`R`$ be large positive constants to be fixed later, $`\chi (t)`$ be given by Lemma 4.1, and let $`\psi 𝒞_\mathrm{c}^{\mathrm{}}(;[0,1])`$ be equal to $`1`$ for $`|t|1`$, and to $`0`$ for $`|t|2`$. We define
$$\chi _0(\rho )\stackrel{\mathrm{def}}{=}\frac{\chi (G_0(\rho ))}{G_0(\rho )}\psi \left(\frac{p(\rho )}{\delta }\right)\psi \left(\frac{|x(\rho )|}{R}\right).$$
Then
$$\begin{array}{cc}\hfill H_p(\chi _0G_0)(\rho )& =\chi ^{}(G_0(\rho ))H_pG_0(\rho )\psi \left(\frac{p(\rho )}{\delta }\right)\psi \left(\frac{|x(\rho )|}{R}\right)\hfill \\ & +\frac{1}{R}\chi (G_0(\rho ))\psi \left(\frac{p(\rho )}{\delta }\right)\psi ^{}\left(\frac{|x(\rho )|}{R}\right)H_p(|x|)(\rho ),\hfill \end{array}$$
and
$$H_p(\chi _0G_0)(\rho )C_1\left(\alpha +\frac{T}{R}\right),$$
where $`C_1`$ is independent of $`T`$ and $`R`$: we note that (7.7) guarantees the boundedness of $`H_pG_0`$, and the assumptions on $`p`$ imply that $`H_p(|x|)`$ is uniformly bounded for $`|p|2\delta `$. For any $`\alpha >0`$ we can choose $`T=T(\alpha )`$ such that $`|G_0(\rho )|\alpha T`$ for $`|x(\rho )|3R_0`$, $`|p(\rho )|2\delta `$. We then choose $`\alpha `$ and $`R`$ so that
$$C_0C_1(\alpha +T(\alpha )/R)<\delta _0.$$
Hence for $`|x(\rho )|R_0`$
$$H_pG=C_0\mathrm{log}(1/ϵ)H_p(\chi _0G_0)\delta _0\mathrm{log}(1/ϵ),$$
which is the last statement in (7.22).
It remains to show (7.23) and for simplicity of presentation we replace $`T^{}X`$ with $`^{2n}`$. We first prove that
(7.24)
$$\frac{\widehat{\phi }_\pm (\rho )+Mϵ}{\widehat{\phi }_\pm (\mu )+Mϵ}C_1\frac{\rho \mu }{\sqrt{ϵ}}^2,M0,$$
with constants depending on $`M`$. We can replace $`\widehat{\phi }_\pm +Mϵ`$ with $`\widehat{\phi }_\pm `$, as $`\widehat{\phi }_\pm +Mϵ_M\widehat{\phi }_\pm `$. Thus we claim that,
$$\frac{\widehat{\phi }_\pm (\rho )}{\widehat{\phi }_\pm (\mu )}C_1\frac{\rho \mu }{\sqrt{ϵ}}^2.$$
Since $`\widehat{\phi }_\pm d(,\mathrm{\Gamma }_\pm )^2+ϵ`$, $`\widehat{\phi }_\pm ϵ`$, we have
$$\begin{array}{cc}\hfill \widehat{\phi }_\pm (\rho )& C(d(\rho ,\mathrm{\Gamma }_\pm )^2+ϵ)C(d(\mu ,\mathrm{\Gamma }_\pm )^2+|\mu \rho |^2+ϵ)\hfill \\ & C^{}(\widehat{\phi }_\pm (\mu )+|\mu \rho |^2)=C^{}(\widehat{\phi }_\pm (\mu )+ϵ(\rho \mu )/\sqrt{ϵ}^2)\hfill \\ & 2C^{}\widehat{\phi }_\pm (\mu )(\rho \mu )/\sqrt{ϵ}^2.\hfill \end{array}$$
In the notation of Lemma 7.6, (7.24) gives
$$|\widehat{G}(\rho )\widehat{G}(\mu )|C+2\mathrm{log}(\rho \mu )/\sqrt{ϵ},$$
and with $`\widehat{\chi }𝒞_\mathrm{c}^{\mathrm{}}`$,
$$|\widehat{\chi }(\rho )\widehat{G}(\rho )\widehat{\chi }(\mu )\widehat{G}(\mu )|C|\rho \mu |\mathrm{log}(1/ϵ)+C\mathrm{log}(\rho \mu )/\sqrt{ϵ}.$$
Clearly,
$$|\chi _0(\rho )G_0(\rho )\chi _0(\mu )G_0(\mu )|C|\rho \mu |\mathrm{log}(1/ϵ),$$
and hence to obtain (7.23) we need
$$|\rho \mu |\mathrm{log}(1/ϵ)C\mathrm{log}(\rho \mu )/\sqrt{ϵ}+C,\rho ,\mu Q^{2n}.$$
If we put $`\delta =\sqrt{ϵ}`$, $`t=|\rho \mu |/(C\delta )`$ this becomes
$$\delta \mathrm{log}\frac{1}{\delta }\frac{\mathrm{log}t+1}{t},0t\frac{1}{\delta },$$
and that is clear as $`t(\mathrm{log}t+1)/t`$ is decreasing. ∎
## 8. Proof of the main result
Let $`G`$ be the escape function given in Proposition 7.7, $`ϵ=h/\stackrel{~}{h}`$. and let $`G^w`$ be its Weyl quantization,
$$G^w=𝒪(\mathrm{log}(\stackrel{~}{h}/h)):L^2(X)L^2(X).$$
We define a family of conjugated operators:
(8.1)
$$P_{\theta ,t}\stackrel{\mathrm{def}}{=}e^{tG^w}P_\theta e^{tG^w},\theta =C_0h\mathrm{log}(1/h).$$
It is easy to see that, in the notation of §3.3,
(8.2)
$$\mathrm{exp}(tG^w)\mathrm{\Psi }_{\frac{1}{2}}^{|t|C,0,0}(X),$$
that is $`\mathrm{exp}(tG^w)=B_t^w`$, $`^\alpha B_t=𝒪(h^{|t|C|\alpha |/2}\stackrel{~}{h}^{|\alpha |/2})`$. Finer estimates are however possible thanks to the results of Bony-Chemin . The first of these is given in
###### Lemma 8.1.
Suppose that $`Q\mathrm{\Psi }_{\frac{1}{2}}^{0,0,0}(^n)`$. Then
(8.3)
$$\mathrm{exp}(tG^w)Q\mathrm{exp}(tG^w)\mathrm{\Psi }_{\frac{1}{2}}^{0,0,0}(^n).$$
###### Proof.
We follow §3.3 and change to the variables
$$(\stackrel{~}{x},\stackrel{~}{\xi })=(\stackrel{~}{h}/h)^{\frac{1}{2}}(x,\xi ),$$
$$\stackrel{~}{G}(\stackrel{~}{x},\stackrel{~}{\xi })=G(x,\xi ),\stackrel{~}{Q}_t(\stackrel{~}{x},\stackrel{~}{\xi })=Q_t(x,\xi ),$$
$$U^1G^w(x,hD)U=\stackrel{~}{G}^w(\stackrel{~}{x},\stackrel{~}{h}D_{\stackrel{~}{x}}),U^1Q_t^w(x,hD)U=\stackrel{~}{Q}_t^w(\stackrel{~}{x},\stackrel{~}{h}D_{\stackrel{~}{x}}),$$
$$Uv(\stackrel{~}{x})=(\stackrel{~}{h}/h)^{\frac{n}{4}}v((h/\stackrel{~}{h})^{\frac{1}{2}}\stackrel{~}{x}).$$
We also note that
$$R\mathrm{\Psi }_{\frac{1}{2}}^{0,0,0}(^n)U^1RU\mathrm{\Psi }^{0,0}(^n),$$
where on the right, $`\stackrel{~}{h}`$ is the small parameter – see the proof of Lemma 3.5. The estimate (7.23) shows that, in $`(\stackrel{~}{x},\stackrel{~}{\xi })`$ coordinates, $`\stackrel{~}{G}`$ satisfies the hypothesis of Proposition 3.7 and that proves (8.3). ∎
The basic properties of $`P_{t,\theta }`$ are given in
###### Proposition 8.2.
Let $`P_{\theta ,t}`$ be given by (8.1) and let $`\mathrm{\Sigma }T^{}X`$ be a compact surface coinciding with $`p^1(0)`$ in a neighbourhood of the support of $`G`$. Then for $`|t|C`$,
(8.6)
$$\begin{array}{c}P_{\theta ,t}=P_\theta ith\mathrm{Op}_h^w(H_pG)+E_t,P_\theta ith\mathrm{Op}_h^w(H_pG)\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^{0,0,2,0}(X)\mathrm{\Psi }_{\frac{1}{2}}^{0,0,2}(X),\\ E_t\mathrm{\Psi }_{\frac{1}{2}}^{1,1,0}(X),E_t=𝒪(h\stackrel{~}{h}):L^2(X)L^2(X),\end{array}$$
uniformly in $`h`$ and $`\stackrel{~}{h}`$.
###### Proof.
Let $`V_1,V_2`$ be open neighbourhoods of $`\mathrm{supp}G`$,
$$\mathrm{supp}GV_1\overline{V}_2T^{}X.$$
We first observe that if $`\mathrm{\Psi }\mathrm{\Psi }^{0,\mathrm{}}(X)`$ satisfies
$$\mathrm{WF}_h(\mathrm{\Psi })V_2,\mathrm{WF}_h(I\mathrm{\Psi })\mathrm{}V_1,$$
then
(8.7)
$$[\mathrm{exp}(tG^w),\mathrm{\Psi }]\mathrm{\Psi }^\mathrm{},\mathrm{}(X),(I\mathrm{\Psi })(\mathrm{exp}(tG^w)I)\mathrm{\Psi }^\mathrm{},\mathrm{}(X),|t|1.$$
In fact, using the calculus in §3.3 we see that $`[G^w,\mathrm{\Psi }]\mathrm{\Psi }^\mathrm{},\mathrm{}(X)`$, Hence, using (8.2)
$$\begin{array}{cc}\hfill \frac{d}{dt}[\mathrm{exp}(tG^w),\mathrm{\Psi }]& =G^w[\mathrm{exp}(tG^w),\mathrm{\Psi }]+[G^w,\mathrm{\Psi }]\mathrm{exp}(tG^w)\hfill \\ & =G^w[\mathrm{exp}(tG^w),\mathrm{\Psi }]+A_t,A_t\mathrm{\Psi }^\mathrm{},\mathrm{}(X).\hfill \end{array}$$
Thus
$$[\mathrm{exp}(tG^w),\mathrm{\Psi }]=_0^t\mathrm{exp}((ts)G^w)A_s𝑑s\mathrm{\Psi }^\mathrm{},\mathrm{}(X),$$
which is the first statement in (8.7). We also compute
$$\frac{d}{dt}(I\mathrm{\Psi })(\mathrm{exp}(tG^w)I)=(I\mathrm{\Psi })G^w\mathrm{exp}(tG^w)\mathrm{\Psi }^\mathrm{},\mathrm{}(X),$$
and the second statement in (8.7) follows. Treating the equivalence of $`(I\mathrm{\Psi })P_\theta e^{tG^w}`$ and $`(I\mathrm{\Psi })P_\theta `$ similarly we conclude that
$$P_{\theta ,t}e^{tG^w}\mathrm{\Psi }P_\theta e^{tG^w}(I\mathrm{\Psi })P_\theta \mathrm{\Psi }^\mathrm{},\mathrm{}(X).$$
We now put
$$Q_\theta \stackrel{\mathrm{def}}{=}\mathrm{\Psi }P_\theta \mathrm{\Psi }^{0,0}(X),Q_{\theta ,t}\stackrel{\mathrm{def}}{=}e^{tG^w}Q_\theta e^{tG^w},$$
and we only need to prove (8.6) with $`P_{}`$ replaced by $`Q_{}`$. By a localization argument similar to the one used to construct $`Q_{}`$ we can assume that $`X=^n`$ when applying Lemma 8.1 and that shows that
$$Q_{\theta ,t}\mathrm{\Psi }_{\frac{1}{2}}^{0,0,0}(X).$$
We now establish the expansion in (8.6). Lemma 3.6 implies that
$$[Q_\theta ,G^w]=(h/i)\mathrm{Op}_h^w(H_{p_\theta }G)+R,$$
where $`R\mathrm{\Psi }_{\frac{1}{2}}^{3/2,3/2,0}(X)\mathrm{\Psi }_{\frac{1}{2}}^{1,1,0}(X)`$. It also shows that
$$[[Q_\theta ,G^w],G^w]=(h/i)[\mathrm{Op}_h^w(H_{p_\theta }G),G^w]+[R,G^w]\mathrm{\Psi }_{\frac{1}{2}}^{1,1,0}(X).$$
Here we used the special structure of $`G`$,
$$G=\widehat{\chi }\widehat{G}+C_0\mathrm{log}(1/h)\chi _0G_0,$$
where $`\widehat{\chi },\chi _0`$ and $`G_0`$ are uniformly smooth. When derivatives fall on these terms in error estimates (3.22) the gain in $`h`$ compensates for the logarithmic growth, while for $`|\alpha |>0`$, $`^\alpha \widehat{G}S_{\frac{1}{2}}^{|\alpha |/2,|\alpha |/2}`$.
This gives,
$$\frac{d}{dt}E_t=[Q_{\theta ,t},G^w](h/i)\mathrm{Op}_h^w(H_{p_\theta }G)+(h/i)\mathrm{Op}_h^w(H_{p_\theta p}G)=[Q_{\theta ,t}Q_\theta ,G^w]+R_t,$$
with
$$E_0=(h/i)\mathrm{Op}_h^w(H_{p_\theta p}G)(h\mathrm{log}(1/h))^2\mathrm{\Psi }_{\frac{1}{2}}^{0,0,0}(X)\mathrm{\Psi }_{\frac{1}{2}}^{1,1,0}(X),$$
and $`R_t\mathrm{\Psi }_{\frac{1}{2}}^{1,1,0}(X)`$. We also have
$$\frac{d}{dt}[(Q_{\theta ,t}Q_\theta ),G^w]=e^{tG^w}[[Q_\theta ,G^w],G^w]e^{tG^w}\mathrm{\Psi }_{\frac{1}{2}}^{1,1,0}(X),Q_{\theta ,0}Q_\theta =0.$$
Hence $`[Q_{\theta ,t}Q_\theta ,G^w]\mathrm{\Psi }_{\frac{1}{2}}^{1,1,0}(X)`$, and consequently $`E_t\mathrm{\Psi }_{\frac{1}{2}}^{1,1,0}`$.
To show that
$$Q_{\theta ,t}^0\stackrel{\mathrm{def}}{=}Q_\theta thi\mathrm{Op}_h^w(H_pG)\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^{0,0,0,0}(X),$$
it suffices to show, in view of Definition 2, that for any sequence
$$\{a_j\}_{j=1}^MS^0(T^{}X),$$
we have
$$\mathrm{ad}_{\mathrm{Op}_h^w(a_1)}\mathrm{}\mathrm{ad}_{\mathrm{Op}_h^w(a_M)}\mathrm{ad}_P^kQ_{\theta ,t}^0u_{L^2(X)}Ch^{M/2+k}\stackrel{~}{h}^{M/2}u_{L^2(X)}.$$
This will follow if we show that
$$\mathrm{ad}_P^kQ_{\theta ,t}^0\mathrm{\Psi }_{\frac{1}{2}}^{k,0,0}(X),$$
and since that is clear for $`Q_\theta `$ we only need to check that
(8.8)
$$h\mathrm{ad}_P^k\mathrm{Op}_h^w(H_pG)\mathrm{\Psi }_{\frac{1}{2}}^{k,0,0}(X).$$
This follows from the following stronger result
###### Lemma 8.3.
Let $`G`$ and $`P`$ be as above. Then for any $`ϵ>0`$
$$\mathrm{ad}_P^{\mathrm{}}(H_pG)^w\mathrm{\Psi }_{\frac{1}{2}}^{\mathrm{}+ϵ,0,0}(X),\mathrm{},\mathrm{}1.$$
###### Proof.
We start by proving that
(8.9)
$$\mathrm{ad}_P^{\mathrm{}}(H_pG_1)^w\mathrm{\Psi }_{\frac{1}{2}}^{\mathrm{}+ϵ,0,0}(X),\mathrm{}1,$$
where $`G_1=\widehat{G}\widehat{\chi }`$ is given in Proposition 7.7. We claim that $`\mathrm{ad}_P^{\mathrm{}}(H_pG_1)^w=E_{\mathrm{}}^w`$, where
$$E_{\mathrm{}}\left(\frac{h}{i}\right)^{\mathrm{}}\left(H_p^{\mathrm{}+1}G_1+\underset{s=0}{\overset{\mathrm{}1}{}}\underset{r=2(\mathrm{}s)}{\overset{\mathrm{}}{}}h^rV_{rs}^{\mathrm{}}(H_p^{s+1}G_1)\right),$$
with $`V_{rs}^{\mathrm{}}`$ a differential operator of order less than or equal to $`r+\mathrm{}s`$, and where in view of good symbolic properties of $`P`$ the error is $`𝒪(h^{\mathrm{}})`$. In fact, Lemma 3.6 gives for any $`aS_{\frac{1}{2}}^{0,0,0}`$, $`\mathrm{ad}_Pa^w=a_1^w`$,
$$a_1(h/i)\left(H_pa+\underset{r=2}{\overset{\mathrm{}}{}}h^rV_r(a)\right),$$
where $`V_r`$ is a differential operator of order $`r+1`$. The expansion for $`E_{\mathrm{}}`$ comes from iterating this and observing that for any differential operator $`B`$ of order $`q`$
(8.10)
$$H_p^mBH_p^kG_1=\underset{s=0}{\overset{m}{}}B_sH_p^{k+s}G_1,$$
where each $`B_s`$ is a differential operator of order $`q`$. Using (7.15) we now see that for $`r2(\mathrm{}s)`$,
$$h^{\mathrm{}}h^rV_{rs}^{\mathrm{}}(H_p^{s+1}G_1)h^{\mathrm{}+r}S_{\frac{1}{2}}^{(r+\mathrm{}s)/2+ϵ,(r+\mathrm{}s)/2,0}S_{\frac{1}{2}}^{\mathrm{}r/2+ϵ+(\mathrm{}s)/2,0,0}S_{\frac{1}{2}}^{\mathrm{}1/2+ϵ,0,0},$$
where the $`ϵ>0`$ correction in the order comes from the term $`(H_p^{s+1}\widehat{\chi })\widehat{G}S_{\frac{1}{2}}^{ϵ,0,0}`$. Hence
$$E_{\mathrm{}}=(h/i)^{\mathrm{}}(H_p^{\mathrm{}+1}G_1)^w+𝒪(h^{\mathrm{}+\frac{1}{2}})\mathrm{\Psi }_{\frac{1}{2}}^{ϵ,0,0}(X)\mathrm{\Psi }_{\frac{1}{2}}^{\mathrm{}+ϵ,0,0}(X),$$
which gives (8.9). On the other hand, again in the notation of Proposition 7.7,
$$\mathrm{ad}_P^{\mathrm{}}(H_pG_0)^w\mathrm{\Psi }^{\mathrm{},0}(X),$$
and consequently
$$\mathrm{ad}_P^{\mathrm{}}(\mathrm{log}(1/h)H_pG_0)^w\mathrm{log}(1/h)\mathrm{\Psi }_{\frac{1}{2}}^{\mathrm{},0,0}(X)\mathrm{\Psi }_{\frac{1}{2}}^{\mathrm{}+ϵ,0,0}(X).$$
Since $`G=G_1+\mathrm{log}(\stackrel{~}{h}/h)\chi _0G_0`$ the lemma follows. ∎
The lemma immediately gives (8.8) completing the proof of Proposition 8.2. ∎
The next lemma follows from Proposition 7.7:
###### Lemma 8.4.
Let $`\widehat{G}`$ be given in Lemma 7.6, $`\chi 𝒞_\mathrm{c}^{\mathrm{}}()`$, and let $`\psi 𝒞_\mathrm{c}^{\mathrm{}}(T^{}X)`$ be one in a fixed small neighbourhood of $`K`$ and zero outside of another sufficiently small neighbourhood of $`K`$. Then
$$\psi (x,\xi )\chi (H_p\widehat{G}(x,\xi ))S_{\mathrm{\Sigma },\frac{1}{2}}^{0,0,\mathrm{},0}(T^{}X).$$
###### Proof.
Lemma 7.6 gives a stronger condition
$$H_p^{\mathrm{}}_\rho ^\alpha H_p^k(\psi (\rho )\chi (H_pG(\rho ))=𝒪((\stackrel{~}{h}/h)^{|\alpha |/2}),$$
as can be verified using (8.10) ∎
As in §6.2 we modify our operator to obtain global invertibility. Thus we define $`aS_{\mathrm{\Sigma },\frac{1}{2}}^{0,0,\mathrm{},\mathrm{}}(T^{}X)`$ as follows
$$a(x,\xi )\stackrel{\mathrm{def}}{=}\chi \left(\frac{\stackrel{~}{h}}{h}p(x,\xi )\right)\chi (H_pG(x,\xi ))\psi (x,\xi ),$$
$$\chi 𝒞_\mathrm{c}^{\mathrm{}}(;[0,1]),\chi (t)1,|t|1,$$
and $`\psi `$ is as in Lemma 8.4. In particular by taking its support close to $`K`$ we can replace $`G`$ by $`\widehat{G}`$ in the definition of $`a`$.
We then put
(8.11)
$$\stackrel{~}{P}_{\theta ,t}=P_{\theta ,t}i(h/\stackrel{~}{h})\stackrel{~}{\mathrm{Op}}_{h,\stackrel{~}{h}}(a)\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^{0,0,2,1}(X),$$
and first treat the region away from the trapped set:
###### Lemma 8.5.
Suppose that $`P_{\theta ,t}`$ is given by (8.11) and $`\mathrm{\Psi }_0\mathrm{\Psi }^{0,0}(T^{}X)`$ satisfies
$$\mathrm{WF}_h(\mathrm{\Psi }_0)K=\mathrm{}.$$
Then for $`u𝒞_\mathrm{c}^{\mathrm{}}(X)`$, $`zD(0,Ch)`$ we have
$$(\stackrel{~}{P}_{\theta ,t}z)\mathrm{\Psi }_0u_{L^2}th\mathrm{\Psi }_0u_{L^2(X)}/C𝒪(h^{\mathrm{}})u_{L^2(X)},$$
$$0<hh_0(\stackrel{~}{h}),0<\stackrel{~}{h}\stackrel{~}{h}_0(t).$$
###### Proof.
Let us assume that $`u=1`$. Once $`h`$ is small enough, $`a0`$ in a neighbourhood of $`\mathrm{WF}_h(\mathrm{\Psi }_0)`$ and Theorem 4 gives
$$(\stackrel{~}{P}_{\theta ,t}z)\mathrm{\Psi }_0u_{L^2}=(P_{\theta ,t}z)\mathrm{\Psi }_0u_{L^2}+𝒪(h\stackrel{~}{h}^{\mathrm{}}).$$
We then observe that $`P_{\theta ,t}\mathrm{\Psi }_{1/2}^{0,0,2}(X)`$, and consequently we can use the simpler calculus of §3.3. Microlocally near $`\mathrm{WF}_h(\mathrm{\Psi }_0)`$, for $`zD(0,Ch)`$, and for $`t`$ sufficiently large, Proposition 7.7 and the choice of the angle of scaling give,
$$P_{\theta ,t}z=\mathrm{Op}_h^w(Rep_\theta Rez)+i\mathrm{Op}_h^w(Imp_\theta ihtH_pGImz)+𝒪_t(h\stackrel{~}{h}+h^2\mathrm{log}(1/h)),$$
$$|Rep_\theta Rez|<\delta Imp_\theta +htH_pG+Imzth/C.$$
(This is the analogue of (4.2) in the non-trapping case of §4.) Lemma 3.3 applied with $`\mathrm{\Psi }_j`$’s such that $`|Rep_\theta Rez|>\delta `$ on $`\mathrm{WF}_h(\mathrm{\Psi }_1)`$ (with $`\mathrm{\Psi }_j`$’s constructed using Lemma 3.2) completes the proof. ∎
Near the trapped set we use the second microlocal calculus to obtain
###### Lemma 8.6.
Suppose that $`P_{\theta ,t}`$ is given by (8.11) and let $`zD(0,Ch)`$. For $`u𝒞_\mathrm{c}^{\mathrm{}}(X)`$, $`u=1`$, with $`\mathrm{WF}_h(u)`$ in a fixed small neighbourhood of $`K`$ we have
(8.12)
$$\begin{array}{cc}& (\stackrel{~}{P}_{\theta ,t}z)u_{L^2(X)}thu_{L^2(X)}/C,0<hh_0(\stackrel{~}{h}),0<\stackrel{~}{h}\stackrel{~}{h}_0(t).\hfill \end{array}$$
provided that $`t`$ is large enough.
###### Proof.
In a small neighbourhood of $`K`$ the operator is microlocally equal to
$$P_t^{\mathrm{}}\stackrel{\mathrm{def}}{=}Pith\stackrel{~}{\mathrm{Op}}_{h,\stackrel{~}{h}}(H_p\widehat{G})i(h/\stackrel{~}{h})\stackrel{~}{\mathrm{Op}}_{h,\stackrel{~}{h}}(a)+𝒪_{L^2L^2}(h\stackrel{~}{h}),$$
that is,
$$(\stackrel{~}{P}_{\theta ,t}z)u_{L^2(X)}=(P_t^{\mathrm{}}z)u_{L^2(X)}+𝒪(h^{\mathrm{}}),u_{L^2(X)}=1,$$
for $`u`$ with $`\mathrm{WF}_h(u)`$ near $`K`$. For $`zD(0,Ch)`$,
$$P_t^{\mathrm{}}Z\stackrel{\mathrm{def}}{=}(\stackrel{~}{h}/h)(P_t^{\mathrm{}}z)\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^{0,0,2,1}(X),Z\stackrel{\mathrm{def}}{=}(\stackrel{~}{h}/h)z,$$
has the symbol given by
$$p^{\mathrm{}}Z=\lambda Zit\stackrel{~}{h}H_pGi\chi (\lambda )\chi (H_pG)+𝒪(\stackrel{~}{h}^2),ZD(0,C\stackrel{~}{h}),\lambda =(\stackrel{~}{h}/h)p.$$
Now let $`\psi _0,\psi _1𝒞_\mathrm{b}^{\mathrm{}}()`$ satisfy
$$\psi _0(t)^2+\psi _1^2(t)=1,\mathrm{supp}\psi _0\{t:\chi (t)=1\},\psi _1(t)0,|t|1/2.$$
As in Lemma 3.2 we can now find two operators $`\mathrm{\Psi }_j^{\mathrm{}}\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^{0,0,0,0}(T^{}X)`$ such that
$$\sigma _{h,\stackrel{~}{h}}(\mathrm{\Psi }_j^{\mathrm{}})=\psi _j(H_pG),(\mathrm{\Psi }_0^{\mathrm{}})^2+(\mathrm{\Psi }_1^{\mathrm{}})^2=Id+𝒪_{L^2L^2}(\stackrel{~}{h}^{\mathrm{}}).$$
In a neighbourhoood of the support of $`\psi _1(H_pG)`$ the operator $`P_\theta ^{\mathrm{}}Z`$ is elliptic in $`\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}^{0,0,2,1}`$:
$$|\lambda Zi\stackrel{~}{h}H_pGi\chi (\lambda )\chi (H_pG)|(|\lambda Z|+\chi (\lambda )\chi (H_pG))/2\lambda /C,$$
when $`ZD(0,C\stackrel{~}{h})`$ and $`\chi (H_pG)>1/2`$, say. This implies (see Lemma 3.1) that for $`u`$ with $`\mathrm{WF}_h(u)`$ near $`K`$, $`u_{L^2(X)}=1`$,
$$(P_t^{\mathrm{}}Z)\mathrm{\Psi }_0^{\mathrm{}}u_{L^2(X)}\mathrm{\Psi }_0^{\mathrm{}}u_{L^2(X)}/C𝒪(\stackrel{~}{h}^{\mathrm{}}),0<\stackrel{~}{h}\stackrel{~}{h}_0.$$
To estimate $`(P_t^{\mathrm{}}Z)\mathrm{\Psi }_1^{\mathrm{}}u`$ from below we proceed as in §4.1. Let
$$B_t^{\mathrm{}}=\frac{1}{2i}\left(P_t^{\mathrm{}}(P_t^{\mathrm{}})^{}\right),$$
so that the Weyl symbol (in the sense of $`\mathrm{\Psi }_{\mathrm{\Sigma },\frac{1}{2}}`$) of $`B_t^{\mathrm{}}`$ is equal to
$$\chi (\lambda )\chi (H_pG)+𝒪(\stackrel{~}{h})t\stackrel{~}{h}H_pG+𝒪_t(\stackrel{~}{h}^2),$$
where we indicated the dependence on $`t`$ in the second bound. Since $`H_pG1/C`$ in a neighbourhood of the support of $`\psi _1(H_p)`$ we see that for $`u`$ with $`\mathrm{WF}_h(u)`$ near $`K`$, $`u_{L^2(X)}`$,
$$B_t^{\mathrm{}}\mathrm{\Psi }_1^{\mathrm{}}u,\mathrm{\Psi }_1^{\mathrm{}}u(t\stackrel{~}{h}𝒪(\stackrel{~}{h})𝒪_t(\stackrel{~}{h}^2))\mathrm{\Psi }_1^{\mathrm{}}u^2𝒪(\stackrel{~}{h}^{\mathrm{}}).$$
We now first take $`t`$ large enough to dominate the first error term and then $`\stackrel{~}{h}`$ small enough to dominate the second one. Hence,
$$\begin{array}{cc}\hfill (P_t^{\mathrm{}}Z)\mathrm{\Psi }_1^{\mathrm{}}u\mathrm{\Psi }_1^{\mathrm{}}u& |(P_t^{\mathrm{}}Z)\mathrm{\Psi }_1^{\mathrm{}}u,\mathrm{\Psi }_1^{\mathrm{}}u||Im(P_t^{\mathrm{}}Z)\mathrm{\Psi }_1^{\mathrm{}}u,\mathrm{\Psi }_1^{\mathrm{}}u|\hfill \\ & =(B_t^{\mathrm{}}ImZ)\mathrm{\Psi }_1^{\mathrm{}}u,\mathrm{\Psi }_1^{\mathrm{}}ut\stackrel{~}{h}\mathrm{\Psi }_1^{\mathrm{}}u^2/2𝒪(\stackrel{~}{h}^{\mathrm{}}),\hfill \end{array}$$
provided that $`t`$ was large enough, and then $`\stackrel{~}{h}`$ small enough. Lemma 3.3 (or rather its proof) gives
$$(P_t^{\mathrm{}}Z)ut\stackrel{~}{h}u/C,tt_01,0<h<h_0(t),$$
for $`u`$ with $`\mathrm{WF}_h(u)`$ near $`K`$.
We complete the proof by writing
$$\begin{array}{cc}\hfill (\stackrel{~}{P}_{\theta ,t}z)u_{L^2(X)}& =(P_t^{\mathrm{}}z)u_{L^2(X)}+𝒪(h^{\mathrm{}})u_{L^2(X)}\hfill \\ & =(h/\stackrel{~}{h})(P_t^{\mathrm{}}Z)u_{L^2(X)}+𝒪(h^{\mathrm{}})u_{L^2(X)}\hfill \\ & thu_{L^2(X)}/C.\hfill \end{array}$$
The two lemmas are now combined using Lemma 3.3 which gives for large $`t`$, $`0<\stackrel{~}{h}\stackrel{~}{h}_0(t)`$, and $`0<h<h_0(t,\stackrel{~}{h})`$, the invertibility of $`\stackrel{~}{P}_{\theta ,t}z`$, $`zD(0,Ch)`$:
$$(\stackrel{~}{P}_{\theta ,t}z)^1=𝒪(1/h):L^2(X)L^2(X).$$
As in §6.2, Theorem 3 is a consequence of writing
(8.13)
$$\stackrel{~}{\mathrm{Op}}_{h,\stackrel{~}{h}}(a)=R+E,\mathrm{rank}(R)=𝒪(h^\nu ),E=𝒪(\stackrel{~}{h}^{\mathrm{}}):L^2(X)L^2(X),$$
$`\nu >\nu (E)`$, where $`m(E)=2\nu (E)+1`$ is the dimension of the trapped set at energy $`E`$, allowing $`\nu =\nu (E)`$ if the trapped set is of pure dimension.
The decomposition (8.13) follows from Proposition 5.10 and the definition of the Minkowski dimension:
$$m_0=2n1sup\{d:\underset{ϵ0}{lim\; sup}ϵ^d\mathrm{vol}(\{\rho p^1(0):d(\rho ,K)<ϵ\})<\mathrm{}\},$$
with the set being of pure dimension if
$$\underset{ϵ0}{lim\; sup}ϵ^{2n+1+m_0}\mathrm{vol}(\{\rho p^1(0):d(\rho ,K)<ϵ\})<\mathrm{}.$$
In other words, for $`ϵ`$ small
$$\mathrm{vol}(\{\rho p^1(0):d(\rho ,K)<ϵ\})Cϵ^{2n1m},m>m_0,$$
and $`m`$ replaceable by $`m_0`$ when $`K`$ is of pure dimension. In particular,
$$\mathrm{vol}(\mathrm{supp}ap^1(0))C_{\stackrel{~}{h}}h^{(2n1m)/2}=C_{\stackrel{~}{h}}h^{n\nu 1},m=2\nu +1>m_0,$$
with equality if $`K`$ is of pure dimension. Since
$$\mathrm{supp}ap^1(0)\underset{\rho K}{}B_\mathrm{\Sigma }(\rho ,M(h/\stackrel{~}{h})^{\frac{1}{2}}),$$
where $`B_\mathrm{\Sigma }`$ are balls in $`\mathrm{\Sigma }`$ with respect to some fixed smooth metric, and since $`K`$ is invariant under the flow, the standard covering arguments (see \[29, Lemma 3.3\]) show that the hypothesis of Proposition 5.10 are satisfied with
$$K(h)C_{\stackrel{~}{h}}h^\nu ,$$
which completes the proof of Theorem 3.
Appendix
We present a direct proof of Proposition 3.7. The hypotheses on $`G`$ in (3.25) are equivalent to the statement that $`\mathrm{exp}(tG)S(m^t)`$, for all $`t`$. We start with
###### Lemma A.1.
Let $`U(t)\stackrel{\mathrm{def}}{=}(\mathrm{exp}tG)^w(x,D):𝒮(^n)𝒮(^n)`$. For $`|t|<ϵ_0(G)`$, the operator $`U(t)`$ is invertible, and
$$U(t)^1=B_t^w(x,D),B_tS(m^t).$$
###### Proof.
We apply the composition formula (3.24) to obtain
$$U(t)U(t)=Id+E_t^w(x,D),E_tS(1).$$
More explicitely we write (see \[10, Proposition 7.7\] and Lemma 3.6 here)
$$\begin{array}{cc}\hfill E_t(x_1,\xi )& =_0^se^{sA(D)}A(D)(e^{tG(x_1,\xi _1)+tG(x_2,\xi _2)})_{x_2=x_1=x,\xi _2=\xi _1=\xi }ds\hfill \\ & =_0^s(it/2)e^{sA(D)}(D_{\xi _1}GD_{x_2}GD_{x_1}GD_{\xi _2}G)e^{tG(x_1,\xi _1)+tG(x_2,\xi _2)}_{x_2=x_1=x,\xi _2=\xi _1=\xi }ds,\hfill \end{array}$$
where $`A(D)=i\sigma (D_{x_1},D_{\xi _1};D_{x_2},D_{\xi _2})/2`$.
Hence $`E_t=t\stackrel{~}{E}_t`$ where $`\stackrel{~}{E}_tS(1)`$ uniformly, and thus
$$E_t^w(x,D)=𝒪(t):L^2(^n)L^2(^n).$$
This shows that for $`|t|`$ small enough $`Id+E_t^w(x,D)`$ is invertible, and Beals’s lemma (see for instance \[10, Proposition 8.3\]) gives
$$(Id+E_t^w(x,D))^1=C_t^w(x,D),C_tS(1).$$
Hence $`B_t=C_t\mathrm{\#}\mathrm{exp}(tG(x,\xi ))S(m^t)`$. ∎
We now observe that
(A.3)
$$\begin{array}{c}\frac{d}{dt}\left(U(t)\mathrm{exp}(tG^w(x,D))\right)=V(t)\mathrm{exp}(tG^w(x,D)),\\ V(t)=A_t^w(x,D),A_tS(m^t).\end{array}$$
In fact, we see that
$$\frac{d}{dt}U(t)=(G\mathrm{exp}(tG))^w(x,D),U(t)G^w(x,D)=(\mathrm{exp}(tG)\mathrm{\#}G)^w(x,D).$$
As before, the composition formula (3.24) gives
$$\mathrm{exp}(tG)\mathrm{\#}GG\mathrm{exp}(tG)=$$
$$_0^1\mathrm{exp}(sA(D))A(D)\mathrm{exp}(tG(x^1,\xi ^1)G(x^2,\xi ^2)_{x^1=x^2=x,\xi ^1=\xi ^2=\xi },$$
$$A(D)=i\sigma (D_{x^1},D_{\xi ^1};D_{x^2},D_{\xi ^2})/2.$$
The hypothesis on $`G`$ shows that $`A(D)\mathrm{exp}(tG(x^1,\xi ^1))G(x^2,\xi ^2)`$ is a sum of terms of the form $`a(x^1,\xi ^1)b(x^2,\xi ^2)`$ where $`aS(m^t)`$ and $`bS(1)`$. The continuity of $`\mathrm{exp}(A(D))`$ on the spaces of symbols (see \[10, Proposition 7.6\]) gives (A.3).
If we put
$$C(t)\stackrel{\mathrm{def}}{=}V(t)U(t)^1,$$
then by Lemma A.1, $`C(t)=c_t^w`$ where $`c_tS(1)`$. Symbolic calculus shows that $`c_t`$ depends smoothly on $`t`$ and
$$(_t+C(t))(U(t)\mathrm{exp}(tG^w(x,D)))=0.$$
The proof of Proposition 3.7 is now reduced to showing
###### Lemma A.2.
Suppose that $`C(t)=c_t^w(x,D)`$, where $`c_tS(1)`$, depends continuously on $`t(ϵ_0,ϵ_0)`$. Then the solution of
(A.4)
$$(_t+C(t))Q(t)=0,Q(0)=q^w(x,D),qS(1),$$
is given by $`Q(t)=q_t(x,D)`$, where $`q_tS(1)`$ depends continuously on $`t(ϵ_0,ϵ_0)`$.
###### Proof.
The Picard existence theorem for ODEs shows that $`Q(t)`$ is bounded on $`L^2`$. If $`\mathrm{}_j(x,\xi )`$ are linear functions on $`T^{}^n`$ then
$$\frac{d}{dt}\mathrm{ad}_{\mathrm{}_1(x,D)}\mathrm{}\mathrm{ad}_{\mathrm{}_N(x,D)}Q(t)+\mathrm{ad}_{\mathrm{}_1(x,D)}\mathrm{}\mathrm{ad}_{\mathrm{}_N(x,D)}(C(t)Q(t))=0,$$
$$\mathrm{ad}_{\mathrm{}_1(x,D)}\mathrm{}\mathrm{ad}_{\mathrm{}_N(x,D)}Q(0):L^2(^n)L^2(^n).$$
If we show that for any choice of $`\mathrm{}_j^{}s`$ and any $`N`$
(A.5)
$$\mathrm{ad}_{\mathrm{}_1(x,D)}\mathrm{}\mathrm{ad}_{\mathrm{}_N(x,D)}Q(t):L^2(^n)L^2(^n),$$
then Beals’s lemma (see \[10, Chapter 8\]) concludes the proof. We proceed by induction on $`N`$:
$$\mathrm{ad}_{\mathrm{}_1(x,D)}\mathrm{}\mathrm{ad}_{\mathrm{}_N(x,D)}(C(t)Q(t))=C(t)\mathrm{ad}_{\mathrm{}_1(x,D)}\mathrm{}\mathrm{ad}_{\mathrm{}_N(x,D)}Q(t)+R(t),$$
where $`R(t)`$ is the sum of terms of the form
$$A_k(t)\mathrm{ad}_{\mathrm{}_1(x,D)}\mathrm{ad}_{\mathrm{}_k(x,D)}Q(t),k<N,A_k(t)=a_k(t)^w,$$
where $`a_k(t)S(1)`$ depend continuously on $`t`$ (this statement can also be proved by induction using the derivation property of $`ad_{\mathrm{}}`$: $`\mathrm{ad}_{\mathrm{}}(CD)=(\mathrm{ad}_{\mathrm{}}C)D+C(\mathrm{ad}_{\mathrm{}}D)`$). Hence by the induction hypothesis $`R(t)`$ is bounded on $`L^2`$, and depends continuously on $`t`$. Thus
$$\left(\frac{d}{dt}+C(t)\right)\mathrm{ad}_{\mathrm{}_1(x,D)}\mathrm{}\mathrm{ad}_{\mathrm{}_N(x,D)}Q(t)=R(t):L^2(^n)L^2(^n).$$
Since (A.5) is valid at $`t=0`$ we obtain it for all $`t(ϵ_0,ϵ_0)`$. ∎
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# Gravitational scattering of massless scalars in QFT and superstring theory
## 1 Introduction
The goal of this paper is to demonstrate an example of how gravity is incorporated in the superstring theory (SST). In particular, we consider gravitational interaction between two massless scalar particles. In quantum theory the quantity reflecting this interaction is the amplitude for a scattering process of two scalar particles by a graviton. First, we can perform this calculation in traditional quantum field theory. Even though gravity can not be fully quantized in QFT, tree-level calculations can be done to give results in agreement with classical theory. Second, we can use SST to get the amplitude for the corresponding process. Here the particles corresponding to the massless scalars and gravitons are massless closed strings of spin-0 and spin-2, respectively. With appropriate identification of coupling constants the results of QFT and SST do agree, which confirms that SST includes gravitation. In particular, it shows an example of how particles in SST having kinematic properties of a graviton (massless, spin-2) also couple as gravitons.
This result is by no means original, and, actually, the argument of existance of gravity in SST can be made much more general . Also various results in this paper appear in other sources and are cited. Therefore, this particular calculation is more of a way to go through some important topics in QFT and SST, rather than an original work.
## 2 Amplitude in QFT
In this section we perform the calculation of the gravitational scattering of two massless scalars in quantum field theory. We only consider the first-order tree-level process $`\varphi \varphi \varphi \varphi `$, that is, with only one intermediate graviton. One Feynman diagram representing such interaction is Fig. 1 (we notate a scalar particle with $`\varphi `$ and graviton with $`h`$).
In order to calculate such amplitude we need two results from QFT: graviton propagator and 2-scalar/graviton interaction vertex. We proceed now with deriving those values from the classical actions.
### 2.1 Graviton propagator
The gravitational action is :
$$S_G=\frac{1}{16\pi G}d^Dx\sqrt{g}R=\frac{1}{2\kappa ^2}d^Dx_G,$$
(2.1)
with
$$_G\sqrt{g}R,\kappa \sqrt{2\pi G}.$$
(2.2)
Here $`G`$ is Newton’s constant, $`g`$ is the determinant of $`g_{\mu \nu }`$ and $`R`$ is the Ricci scalar. This particular normalization of $`S_G`$ is needed to get the correct Einstein’s equation when we add conventionally normalized matter action. Note that we are working in arbitrary number of dimensions $`D`$.
To be able to do perturbative calculation (to first order in our case) we assume the gravitational field is weak, that is, the metric is almost flat :
$$g_{\mu \nu }=\eta _{\mu \nu }+\beta h_{\mu \nu },|\beta h_{\mu \nu }|1,$$
(2.3)
note that we use the convention for $`\eta _{\mu \nu }`$, that differs from but agrees with the other references listed:
$$\eta _{\mu \nu }=\text{diag}(,+,+,+,\mathrm{}).$$
(2.4)
By the arbitrary normalization constant $`\beta `$ in (2.3) we can redefine the normalization of $`h`$ \- we will need this to get the correctly normalized graviton states in quantum field theory. In non-quantum general relativity it is usually taken $`\beta =1`$ in this expansion. Note that $`h_{\mu \nu }`$, like $`g_{\mu \nu }`$, is symmetric. We also need an expression for the corresponding inverse-metric:
$$g^{\mu \nu }=\eta ^{\mu \nu }\beta h^{\mu \nu }+O(h^2),$$
(2.5)
where indices on $`h_{\mu \nu }`$ are raised with $`\eta ^{\mu \nu }`$.
Now we Taylor-expand $`_G`$ in (2.2) to the second power in $`h`$. Up to a total derrivative the result is (this differs by a sign from that in )
$$_h=\frac{\beta ^2}{2}\left[h_\mu _\nu h^{\mu \nu }h^{\mu \rho }_\mu _\nu h_{}^{\nu }{}_{\rho }{}^{}+\frac{1}{2}h_{\mu \nu }^2h^{\mu \nu }\frac{1}{2}h^2h\right].$$
(2.6)
Here $`h=h_{}^{\mu }{}_{\mu }{}^{}`$ and $`^2=_\mu ^\mu `$. The total derrivative in $`_G`$ is irrelevant since it gets integrated in (2.1) over all space and we assume the boundary terms to vanish. This quadratic Lagrangian $`_h`$ in (2.6), when quantized, describes a massless particle, namely, the graviton<sup>1</sup><sup>1</sup>1The problem with quantum gravity only appears when we consider higher orders in $`_G`$ expansion and loop Feynman diagrams - the theory turns out to be non-renormalizable.
Now we consider quantizing this gravitational action, but let’s fix the $`\beta `$ in (2.3) before we go further. It turns out that the value that gives correctly normalized states of $`h`$ is :
$$\beta =2\kappa .$$
(2.7)
With this $`\beta `$, the $`\kappa `$’s cancel in $`S_G`$ and we get:
$$S_hd^Dx\left[h_\mu _\nu h^{\mu \nu }h^{\mu \rho }_\mu _\nu h_{}^{\nu }{}_{\rho }{}^{}+\frac{1}{2}h_{\mu \nu }^2h^{\mu \nu }\frac{1}{2}h^2h\right],$$
(2.8)
which is the action that we will quantize.
We proceed with the quantization of $`h`$ field using Feynman path integral approach (this part, especially Faddeev-Popov gauge fixing, is based on Chapter 9 in ). We need to consider the quantity
$$Z=𝒟he^{iS_h[h]},$$
(2.9)
where $`𝒟h`$ integrates over all inequivalent field configurations. Once we have such an expression, we can deduce the propagator for the field directly from the $`S`$ appearing in (2.9).
The problem with $`S_h`$ is that there is a gauge transformation:
$$h_{\mu \nu }^\xi =h_{\mu \nu }+_\mu \xi _\nu +_\nu \xi _\mu ,$$
(2.10)
for any field $`\xi _\mu `$, that leaves $`S_h`$ invariant. Therefore, simple $`𝒟h`$ in (2.9) includes many equivalent field configurations. To get the correct $`Z`$ we need to factor out from (2.9) the integral over gauge transformations, leaving only the integral over physically different configurations. This is done by Faddeev-Popov gauge fixing procedure, which we briefly describe here.
First, we choose a gauge-fixing function $`G(h)`$ <sup>2</sup><sup>2</sup>2The discussion here also works for $`G_\lambda (h)`$ with an index, that is, for a collection of gauge-fixing functions., such that the constraint $`G(h)=0`$ would pick out *one* field configuration from each set of equivalent ones, that is, fix the gauge. Then we can write (2.9) as :
$`Z={\displaystyle 𝒟he^{iS_h[h]}}`$ $`=\mathrm{\Delta }{\displaystyle 𝒟\xi 𝒟he^{iS_h[h]}\delta (G(h))},`$ (2.11)
$`\mathrm{\Delta }`$ $`det\left({\displaystyle \frac{\delta G(h^\xi )}{\delta \xi }}\right).`$ (2.12)
Here $`h^\xi `$ as in (2.10) and $`\mathrm{\Delta }`$ is a functional determinant arising as the Jacobian for the relevant change of variables. We must choose $`G`$ such that the $`\mathrm{\Delta }`$ is independent of $`h`$, thus, just a constant. We see that we did factor out $`𝒟\xi `$, the integral over gauge transformations, and the inner integral is constrained by functional $`\delta `$, allowing only gauge-fixed, thus inequivalent, field configurations.
Now we consider a particular gauge-fixing function:
$$G_\lambda =_\mu h_{}^{\mu }{}_{\lambda }{}^{}\frac{1}{2}_\lambda hw_\lambda ,$$
(2.13)
where $`w_\lambda `$ is any field. Then $`G_\lambda =0`$ is a generalized Lorentz gauge condition (it is Lorentz gauge for $`w_\lambda =0`$). We can not, though, use (2.11) for direct calculations, because it is, due to the $`\delta `$-function, not in the standard form (2.9). We fix that as follows. Since (2.11) is true for any $`G`$, we can integrate the equation with $`G_\lambda `$ over $`w`$ with any normalized-function weighting and cancel the $`\delta `$-function with this integral. In particular, we use:
$$N(\alpha )𝒟w\mathrm{exp}\left(id^Dx\alpha w_\lambda w^\lambda \right)=1,$$
(2.14)
with some real non-zero coefficient $`\alpha `$ and normalization constant $`N(\alpha )`$ to integrate (2.11). Note that this works on RHS because $`\mathrm{\Delta }`$ is $`w`$-independent. We get:
$`Z`$ $`=N{\displaystyle 𝒟w\mathrm{exp}\left(id^Dx\alpha w_\lambda w^\lambda \right)\mathrm{\Delta }𝒟\xi 𝒟he^{iS_h[h]}\delta (F_\lambda (h)w_\lambda )}`$
$`=N\mathrm{\Delta }\left({\displaystyle 𝒟\xi }\right){\displaystyle 𝒟h\mathrm{exp}\left(iS_h[h]+id^Dx\alpha F_\lambda F^\lambda \right)},`$ (2.15)
where
$$F_\lambda (h)_\mu h_{}^{\mu }{}_{\lambda }{}^{}\frac{1}{2}_\lambda h.$$
(2.16)
Thus in effect the result of Faddeev-Popov procedure is adding a new term to the Lagrangian and factoring out an overall (infinite) constant from Z:
$`Z`$ $`=C{\displaystyle 𝒟h\mathrm{exp}\left(iS_h^{}\right)},`$ (2.17)
$`S_h^{}`$ $`={\displaystyle d^Dx\left(2_h+\alpha F_\lambda F^\lambda \right)}.`$ (2.18)
All the factors, including gauge integral, are now contained in an overall factor $`C`$, which is irrelevant for calculating expectation values. What we are left with is the standard form of $`Z`$ with action $`S_h^{}`$, from which we can correctly deduce the graviton propagator.
The calculations must work for any value of $`\alpha `$, but we can set it to simplify our Lagrangian (2.6). The additional term, up to a total derrivative, is:
$$\alpha F_\lambda F^\lambda =\alpha \left(h_\mu _\nu h^{\mu \nu }h^{\mu \rho }_\mu _\nu h_{}^{\nu }{}_{\rho }{}^{}\frac{1}{4}h^2h\right).$$
(2.19)
Choosing $`\alpha =1`$, action (2.18) becomes:
$$S_h^{}=\frac{1}{2}d^Dx\left(h_{\mu \nu }^2h^{\mu \nu }\frac{1}{4}h^2h\right),$$
(2.20)
which we can finally use to get the graviton propagator.
In general, we can deduce the propagator from an action as follows. Suppose we have a real field $`\varphi `$ described by:
$$S=\frac{1}{2}d^Dx\varphi (x)Q\varphi (x),$$
(2.21)
with some differential operator $`Q`$. In case $`\varphi `$ has indices, $`\varphi Q\varphi `$ is matrix multiplication. The propagator $`D(xy)=\varphi (x)\varphi (y)`$ then satisfies:
$$QD(xy)=i\delta (xy)𝐈,$$
(2.22)
where $`𝐈`$ is matrix identity. With $`Q`$ and $`D`$ transformed into momentum space (2.22) becomes:
$$\stackrel{~}{Q}(k)\stackrel{~}{D}(k)=i𝐈.$$
(2.23)
Now we apply these identities to $`S_h^{}`$. We can rewrite (2.20) as:
$`S_h^{}`$ $`={\displaystyle \frac{1}{2}}{\displaystyle d^Dxh_{\mu \nu }Q^{\mu \nu ;\rho \sigma }h_{\rho \sigma }},`$ (2.24)
$`Q^{\mu \nu ;\rho \sigma }`$ $`={\displaystyle \frac{1}{2}}\left(\eta ^{\mu \rho }\eta ^{\nu \sigma }+\eta ^{\mu \sigma }\eta ^{\nu \rho }\eta ^{\mu \nu }\eta ^{\rho \sigma }\right)^2,`$ (2.25)
where we have explicitly symmetrized $`Q`$ in $`\{\mu \nu \}`$ and $`\{\rho \sigma \}`$, since $`h_{\mu \nu }`$ is always symmetric, thus antisymmetric matrix component acting on it is irrelevant. In momentum space this looks as
$$\stackrel{~}{Q}^{\mu \nu ;\rho \sigma }=\frac{k^2}{2}\left(\eta ^{\mu \rho }\eta ^{\nu \sigma }+\eta ^{\mu \sigma }\eta ^{\nu \rho }\eta ^{\mu \nu }\eta ^{\rho \sigma }\right),$$
(2.26)
and the identity in (2.23), since we are working with symmetric matrices, is:
$$𝐈_{\rho \sigma }^{\mu \nu }=\frac{1}{2}\left(\delta _\rho ^\mu \delta _\sigma ^\nu +\delta _\sigma ^\mu \delta _\rho ^\nu \right).$$
(2.27)
Now we can find $`\stackrel{~}{D}`$ satisfying:
$$\stackrel{~}{Q}^{\mu \nu ;\alpha \beta }\stackrel{~}{D}_{\alpha \beta ;\rho \sigma }=i𝐈_{\rho \sigma }^{\mu \nu },$$
(2.28)
to be:
$$\stackrel{~}{D}_{\mu \nu ;\rho \sigma }^{(h)}(k)=\frac{1}{2}\left(\frac{i}{k^2iϵ}\right)\left(\eta _{\mu \rho }\eta _{\nu \sigma }+\eta _{\mu \sigma }\eta _{\nu \rho }\frac{2}{D2}\eta _{\mu \nu }\eta _{\rho \sigma }\right).$$
(2.29)
in agreement with . This is the final expression for momentum-space graviton propagator between polarizations $`(\mu \nu )`$ and $`(\rho \sigma )`$. Term $`iϵ`$ was added in the denominator as usual to give the right behavior of the integral to position space.
### 2.2 Scalar-graviton interaction
Now we proceed with calculating the second component needed for our scattering calculation: the scalar-graviton interaction vertex. Again we start with classical action and quantize it.
First consider a general matter action added to the pure gravitational one (2.1):
$`S`$ $`=S_G+S_M,`$ (2.30)
$`S_M`$ $`={\displaystyle d^Dx\sqrt{g}_M}.`$ (2.31)
This action includes matter-gravity interaction, because it has both metric and matter field terms. Note that the Lagrangian $`_M`$ itself can contain metric terms.
Again we want to consider weak gravitational field behavior, for which we expand $`S_M`$ in Taylor series around the flat matric. As in (2.3), variation in metric is $`\beta h_{\mu \nu }`$, and we will set $`\beta =2\kappa `$ when quantizing. We are interested only in first-order interaction, which corresponds to vertices with only one graviton involved, so we expand $`S_M`$ up to the first order in $`h`$:
$`S_M=(S_M)_{g=\eta }+{\displaystyle d^Dx(\beta h^{\mu \nu })\left(\frac{\delta S_M}{\delta g^{\mu \nu }}\right)_{g=\eta }}+O(h^2).`$ (2.32)
The first term describes non-interacting (gravitationally) matter:
$$S_m(S_M)_{g=\eta }=d^Dx(_M)_{g=\eta }.$$
(2.33)
The second describes first-order gravitational interaction and it can be expressed in terms of a familiar quantity - energy-momentum tensor, defined as (, up to a factor):
$$T_{\mu \nu }\frac{2}{\sqrt{g}}\frac{\delta S_M}{\delta g^{\mu \nu }}=2\frac{_M}{g^{\mu \nu }}+g_{\mu \nu }_M.$$
(2.34)
The interaction term in (2.32) is then:
$$S_I=\frac{\beta }{2}d^Dxh^{\mu \nu }(T_{\mu \nu })_{g=\eta },$$
(2.35)
and the total matter action is:
$$S_M=S_m+S_I+O(h^2).$$
(2.36)
Now we go to our case of interest, the massless scalar field $`\varphi `$. Such field is described by the action:
$`S_M`$ $`={\displaystyle \frac{1}{2}}{\displaystyle d^Dx\sqrt{g}(g^{\mu \nu }_\mu \varphi _\nu \varphi )},`$ (2.37)
$`_M`$ $`={\displaystyle \frac{1}{2}}g^{\mu \nu }_\mu \varphi _\nu \varphi .`$ (2.38)
We get the free field action $`S_m`$ by just replacing $`g^{\mu \nu }`$ by $`\eta ^{\mu \nu }`$ in (2.37), from which we could easily get the scalar field propagator. For our calculation that is not necessary, though, so we concentrate on $`S_I`$, which contains the interaction vertex we need. To use (2.35) first we get calculate the momentum-energy tensor for $`\varphi `$ by (2.34):
$$T_{\mu \nu }=_\mu \varphi _\nu \varphi \frac{1}{2}g_{\mu \nu }(_\lambda \varphi ^\lambda \varphi ).$$
(2.39)
Now plugging this $`T_{\mu \nu }`$ in (2.35) with $`\beta =2\kappa `$ we have:
$$S_I=\kappa d^Dx\left(h_{\mu \nu }^\mu \varphi ^\nu \varphi \frac{1}{2}h_{\mu \nu }\eta ^{\mu \nu }_\lambda \varphi ^\lambda \varphi \right),$$
(2.40)
the action which we use for quantization.
In the Feynman path integral quantization the vertex values can be easily read off directly from the action - they are basically just the coefficients of the corresponding field products. In this $`S_I`$ we have a product of two $`\varphi `$ fields and an $`h_{\mu \nu }`$ field, which does correspond to the vertex we want (Fig. 2).
We arbitrarily chose the momenta to be incoming - the outgoing momentum is then represented by negative-incoming. To get the vertex value in the momentum space we first substitute $`_\mu ik_\mu `$ (for an incoming momentum), with $`k_\mu `$ of the field the derrivative acts on, which gives the term in $`S_I`$ as:
$$\kappa \left(k_1^\mu k_2^\nu \frac{1}{2}\eta ^{\mu \nu }k_1k_2\right)\varphi (k_1)\varphi (k_2)h_{\mu \nu }(k_3).$$
The value of the vertex is then given by this coefficient in front of the fields multiplied by a factor of 2 from permuting identical lines $`\varphi `$ and an a factor of $`i`$ that we have in any vertex (it comes from $`\mathrm{exp}(iS)`$). Thus the vertex amplitude is:
$$V_{(\varphi \varphi h)}^{\mu \nu }(k_1,k_2)=i\kappa \left(k_1^\mu k_2^\nu +k_1^\nu k_2^\mu \eta ^{\mu \nu }(k_1k_2)\right),$$
(2.41)
where we have again explicitly symmetrized the tensor multiplying $`h_{\mu \nu }`$.
### 2.3 QFT scattering amplitude
Now we have all the pieces to calculate the scattering amplitude (S-matrix) for our process. Consider the Feynman diagram Fig. 3.
At first we perform the amplitude calculation only for this one “channel” - the diagram where $`k_1`$ and $`k_3`$ are connected at one vertex ($`t`$-channel), and then we construct the full amplitude (which will include terms from two more channels) using this result. Note that we marked all $`k_i`$ as incoming - this is partly to make the comparison with the string theory result easier. In applying our result based on this diagram to an actual process with initial momenta $`(k_1,k_2)`$ and final momenta $`(k_1^{},k_2^{})`$ we would have to take $`k_3=k_1^{}`$ and $`k_4=k_2^{}`$ to reverse the direction.
The S-matrix for the diagram Fig. 3 is (Chapter 4.5 in ):
$`S_t`$ $`=(2\pi )^D\delta ^D(\mathrm{\Sigma }_ik_i)i_t,`$ (2.42)
$`i_t`$ $`=V_{(\varphi \varphi h)}^{\mu \nu }(k_1,k_2)\stackrel{~}{D}_{\mu \nu ;\rho \sigma }^{(h)}(k_1+k_2)V_{(\varphi \varphi h)}^{\rho \sigma }(k_3,k_4).`$ (2.43)
We can evaluate $`i_t`$ by plugging in values from (2.29) and (2.41). Note that for our case of interest:
$`k_1^2=k_2^2=k_3^2=k_4^2=0,`$ (2.44)
$`k_1+k_2+k_3+k_4=0,`$ (2.45)
due to masslessness and momentum conservation ($`\delta `$-function in (2.42)) respectively. It is useful then to define “Mandelstam variables” :
$`s`$ $`(k_1+k_2)^2=2k_1k_2=(k_3+k_4)^2=2k_3k_4,`$ (2.46)
$`t`$ $`(k_1+k_3)^2=2k_1k_3=(k_2+k_4)^2=2k_2k_4,`$ (2.47)
$`u`$ $`(k_1+k_4)^2=2k_1k_4=(k_2+k_3)^2=2k_2k_3,`$ (2.48)
that also satisfy
$$s+t+u=0.$$
(2.49)
Evaluating $`i_t`$ gives our desired amplitude for $`t`$-channel, which has a very simple expression in terms of Mandelstam variables:
$`i_t`$ $`=i\kappa ^2{\displaystyle \frac{su}{t}},`$ (2.50)
$`S_t`$ $`=i\kappa ^2(2\pi )^D\delta ^D(\mathrm{\Sigma }_ik_i){\displaystyle \frac{su}{t}}.`$ (2.51)
Note that the there is a pole in $`S`$ at $`t=0`$, which should have been expected, since $`t=(k_1+k_3)^2=q^2=m^2`$ of the virtual graviton, and the pole appears when the virtual particle is on-shell, in this case, $`m^2=0`$. It is because of the pole in $`t`$, which comes from the fact that $`k_1`$ and $`k_2`$ are connected at a vertex, that this configuration is called $`t`$-channel.
As mentioned before, the amplitude $`S_t`$ is not the full amplitude because there are two more Feynman diagrams contributing to the same observed scattering process $`\varphi \varphi \varphi \varphi `$ mediated by a graviton. We get these by simply permuting $`k_i`$’s in the diagram of Fig. 3 \- the resulting diagrams are shown in Fig. 4.
The amplitudes are as in (2.43), but with $`k_i`$ permuted, which causes $`s,u,t`$ permutation:
$`i_u`$ $`=i\kappa ^2{\displaystyle \frac{st}{u}},`$ (2.52)
$`i_s`$ $`=i\kappa ^2{\displaystyle \frac{tu}{s}}.`$ (2.53)
These processes have poles in $`u`$ and $`s`$ (from $`(k_1,k_4)`$ and $`(k_1,k_3)`$ connections), therefore are called $`u`$-channel and $`s`$-channel.
The full scattering amplitude for the process is thus the sum of the three diagrams, All the particles involved are bosons, so there are also no relative minus signs from permutations, therefore, the total amplitude is:
$$S=S_t+S_u+S_s=i\kappa ^2(2\pi )^D\delta ^D(\mathrm{\Sigma }_ik_i)\left(\frac{tu}{s}+\frac{us}{t}+\frac{st}{u}\right),$$
(2.54)
which is our final QFT result that we will compare to the string theory calculation.
## 3 Amplitude in SST
Now we proceed with the calculation of the scattering amplitude in superstring theory. This calculation is strongly based on refs. and many of the results in this section appear there.
We consider Type I open and Type II closed superstring theories, which are among the most realistic string theories, and include our particles of interest - massless scalars and massless spin-2 particles (gravitons). For our purposes heterotic superstring theory could be substituted for Type II, as it contains all the same properties we are interested in, and the result for this scattering in comparison with QFT would be the same.
The superstring theories require a 10-dimensional spacetime. We take it to be flat, that is, described by metric $`\eta _{\mu \nu }`$ as in (2.4), with 9 extended spatial dimensions. This does not represent a realistic model for our universe, but considering compactifications of 6 dimensions adds too many complications. We can still compare our result in 10 dimensions with the QFT result, as we derived it for arbitrary $`D`$. We will continue denoting the number of dimensions as $`D`$ in this section, but it should be kept in mind that the results only make sense for
$$D=10.$$
(3.1)
### 3.1 General amplitudes in SST
Here we will discuss qualitatively how the strings are described and, consequently, how the amplitudes are calculated in string theory. The quantitative results needed for the calculation, the string vertex operators and related field expectation values, will not be derived, and taken from in the following sections.
Strings moving in spacetime are described by a collection of *fields* living on a two-dimensional surface called *world sheet*. To see this consider a motion of an open classical string, which can be described by its position in spacetime $`X^\mu (\sigma ,\tau )`$, with *two* parameters: $`\sigma (0,a)`$ parametrizes position along the string and $`\tau (\mathrm{},\mathrm{})`$ is some parametrization of the motion in time. Then we can interpret this parameter space as a two-dimensional space (world sheet), on which we have $`D`$ scalar fields $`X^\mu `$. In the case of the open string, the world sheet is a strip of finite width and infinite length, while for a closed string it would be an infinitely long cylinder because of the identification $`\sigma \sigma +a`$.
The quantization of a string is then just the quantization of the fields with the appropriate action on the world sheet, which results in a two-dimensional QFT. In the quantization, however, more fields than just $`X^\mu `$ are introduced: for a superstring we have in addition $`D`$ anticommuting fields $`\psi ^\mu `$, also the “world-sheet metric” that appears in the action is itself made into a dynamic field and, finally, we get gauge symmetries which result in ghost fields while being fixed. Each of these fields makes up an independent QFT, with the expectation values that we will cite when needed.
One important feature of the string action is a *conformal symmetry* of the world-sheet, which means that the world-sheets related by a conformal transformation are equivalent, that is, they describe the same string process. This symmetry has many important consequences and the QFT’s with such symmetry are called conformal field theories (CFT’s). The one property of conformal symmetry relevant to our qualitative discussion is that we essentially need to consider only different *topologies* of the world sheets, ignoring the exact shape<sup>3</sup><sup>3</sup>3Conformal equivalence of world-sheets is a bit narrower definition than topological equivalence (meaning that one can be deformed continuously into another). The difference, however, is “small” in a sense that the parameter space of conformally-inequivalent world-sheets within a certain topology is finite-dimensional..
Now we consider string interactions. The nice property of string theory is that interactions are described in a natural way by the same acion, only with considering non-trivial shapes of the world-sheets that the fields live on. For example in Fig. 5 we see a world sheet for 4 open strings interacting by an open string and in Fig. 6 for 4 closed strings interacting by another closed string. We can conclude that it is that particular kind of intermediate string by looking at a shape of a cut through the world sheet in the “interaction region”.
Shown with arrows are conformally equivalent world sheets (it is easy to see that they are topologically equivalent). In case of 4 open strings we get a disk with 4 points missing at the boundary - these arise from 4 “legs” extending infinitely in time-like direction of the world sheet, where the boundary is missing. Similarly for 4 closed strings the conformally equivalent surface is a sphere with 4 points missing.
Consider, finally, the calculation of the amplitude itself. By the method of Feynman path integral quantization, it is given by the sum of $`e^{iS}`$ over all inequivalent configurations of world sheets $`𝒲`$ and fields on them $`\mathrm{\Phi }`$, with some specified initial and final states:
$$S=𝒟[𝒲]𝒟[\mathrm{\Phi }]_{\mathrm{\Phi }_0(z_0)}\mathrm{exp}(iS[𝒲,\mathrm{\Phi }]).$$
(3.2)
We denoted by $`\mathrm{\Phi }_0(z_0)`$ the constraint on $`\mathrm{\Phi }`$ that it has to satisfy at the set of points $`z_0`$ ($`𝒲`$-dependent) in the boundary of the world-sheet, where the initial and final string states are defined. We can always conformally transform the world-sheet so that, as in our examples, the initial and final state strings appear as missing *points* $`z_i`$, around which $`\mathrm{\Phi }_0`$ has to be defined. It can be shown that equivalently we can substitute the initial and final state constraints by a string *vertex operators* $`𝒱_i(z_i)`$ at those points, and let $`\mathrm{\Phi }`$ take any values:
$$S=𝒟[𝒲]𝒟[\mathrm{\Phi }]\mathrm{exp}(iS[𝒲,\mathrm{\Phi }])\underset{i}{}𝒱_i(z_i).$$
(3.3)
This says that the S-matrix amplitudes are the expectation values of vertex operators, which we can think of as operators creating an incoming or outgoing string somewhere on the world sheet.
Consider now splitting the integral $`𝒟[𝒲]`$ into a sum of integrals within each topology:
$$S=\underset{\text{topologies}}{}𝒟[]𝒟[\mathrm{\Phi }]\mathrm{exp}(iS[𝒲,\mathrm{\Phi }])\underset{i}{}𝒱_i(z_i).$$
(3.4)
The $``$ is the so-called *moduli* \- the parameters within a topology distinguishing conformally-inequivalent world sheets (this is where the small difference between topological and conformal equivalence goes). We concentrate instead on the sum over topologies now. The world-sheet topology tells us how many and what kind of strings (closed, open) are involved in the process - it is the equivalent of Feynman diagrams in QFT. More complicated topologies (that is, with more holes) will describe processes with more interactions and more intermediate particles. This way we organized the calculation as perturbation series, and we can constrain ourselves to simpler topologies if we are interested in the leading terms in the series.
In particular, the world-sheet topologies in Fig. 5 and Fig. 6 correspond to tree-level diagrams for 4 open and closed strings respectively. Any other topologies would have holes, therefore, would represent loop-diagrams. Since we want to compare our calculation here with a tree-level result in QFT, we will be interested in these tree-level topologies. Note also, that unlike the Feynman diagrams, where we had 3 different ones for different channels, the one topology in Fig. 6, for example, will contain all the 3 possibilities. It follows from the fact that the string interaction is not exactly localized on the world sheet, and we can get different strings merging into intermediate one by just streching the world sheet differently, which will still be the same topology. Therefore, all the possible tree-level interactions of 4 closed strings by a closed string is represented by one topology - a sphere.
### 3.2 String states
We will describe now the string states that we will need for the calculation. In string theory all the different particles are just different excitations of a string, so we need to find what excitations correspond to our primary particles of interest: gravitons and massless scalars. It turns out, that if we look for massless bosons, of spin-2 and spin-0, we find that both correspond to massless closed strings in NS-NS <sup>4</sup><sup>4</sup>4In superstring theory there are two types of string ground states, therefore, the string states, created by excitations of a ground state, fall into two disctinct sets, which are called *sectors*. One is called Neveu-Schwarz (NS) sector and contains bosons. The other one, Ramond (R) sector, contains fermions. Closed strings are basically composed of two such states and, therefore, there are four sectors: NS-NS and R-R are bosons, and NS-R and R-NS are fermions. sector with different polarizations. Massless NS-NS states are constructing by acting on a $`k`$-momentum NS-NS ground state with two vector creation operators :
$$\psi _{1/2}^\mu \stackrel{~}{\psi }_{1/2}^\nu |0;k_{\text{NS-NS}},$$
(3.5)
A general physical massless closed string state is then described by momentum $`k^\mu `$ and polarization 2-tensor $`e_{\mu \nu }`$:
$$|e;k_{\text{NS-NS}}=e_{\mu \nu }\psi _{1/2}^\mu \stackrel{~}{\psi }_{1/2}^\nu |0;k_{\text{NS-NS}},$$
(3.6)
which in addition have to satisfy the requirements :
$$k^2=k^\mu e_{\mu \nu }=k^\nu e_{\mu \nu }=0,$$
(3.7)
and the states are physically identified under:
$$e_{\mu \nu }e_{\mu \nu }+a_\mu k_\nu +k_\mu b_\nu ,ak=bk=0.$$
(3.8)
The reason why these types of excitations represent more than just one particle, is that the 2-tensor representation $`e_{\mu \nu }`$ by which the particles transform under the *little group* $`SO(D2)`$ that leaves the momentum invariant, is not an irreducible representation. It decomposes into a *symmetric traceless tensor*:
$$e_{\mu \nu }^{(g)}=e_{\nu \mu }^{(g)},\eta ^{\mu \nu }e_{\mu \nu }^{(g)}=0,$$
(3.9)
an *antisymmetric tensor*:
$$e_{\mu \nu }^{(b)}=e_{\nu \mu }^{(b)},$$
(3.10)
and a *trace part*, which is invariant under the $`SO(D2)`$ :
$$e_{\mu \nu }^{(\varphi )}=\frac{1}{\sqrt{D2}}\left(\eta _{\mu \nu }k_\mu \xi _\nu \xi _\mu k_\nu \right),k\xi =1.$$
(3.11)
These three components each satisfy the constraints (3.7) and don’t mix under the transformations of the little group. The symmetric traceless part $`e^{(g)}`$ corresponds to the *graviton* and the trace part $`e^{(\varphi )}`$ is $`SO(D2)`$ invariant, thus a *massless scalar* particle, which in string theory is called *dilaton*. Therefore, we will be looking for an amplitude for 4 $`e^{(\varphi )}`$-polarized closed string scattering by a $`e^{(g)}`$-polarized closed string. Note that we chose the normalization for the explicit expression $`e_{\mu \nu }^{(\varphi )}`$ such that
$$e_{\mu \nu }e^{\mu \nu }=1.$$
(3.12)
Such normalization is assumed by the expressions for vertex operators and amplitudes that we will cite.
Even though we don’t necessarily have to deal with open strings in this calculation it will be very convenient to use them for an intermediate step in the calculation. It is because, as we will see, the amplitude for a closed string essentially factors into a product of two open-string amplitudes, which are in turn twice as less complicated. For this calculation we will need to know the states of massless open strings. Again we are interested in the NS sector, as the amplitude of NS-NS closed strings will be composed of two NS open string amplitudes. The general state here is :
$$|e;k_{\text{NS}}=e_\mu \psi _{1/2}^\mu |0;k_{\text{NS}},$$
(3.13)
with
$$k^2=k^\mu e_\mu =0,$$
(3.14)
and identification:
$$e_\mu e_\mu +\gamma k_\mu .$$
(3.15)
The states here transform as a vector $`e_\mu `$, thus are irreducible and represent one particle. This massless vector boson is actually identified with a photon. Again, it will be assumed that $`e_\mu e^\mu =1`$.
### 3.3 Open string tree amplitudes
We proceed now with calculation of the amplitudes for open string scattering - in particular, we will calculate tree-level amplitudes for interactions of 3 and 4 massless bosons described by the string state (3.13). We are primarily interested in the *closed* string scattering, however, the results in this section will allow us to get the closed string amplitudes quickly.
First consider 3 open string interaction. The tree-level world surface for this process is a disk with three vertex operators at the boundary. This either describes 2 strings joining or 1 string splitting into two - the amplitude we calculate is general, because whether the string is incoming or outgoing depends only on the momentum $`k`$ in the vertex operator - for outgoing string $`k`$ is the negative of the actual string momentum.
From the conformally equivalent shapes to represent a disk, for all explicit expressions we fix it to be the upper half of the complex plane, with the boundary being the real line. The vertex operators for the open string state (3.13) inserted at position $`y(\mathrm{},\mathrm{})`$ on the world-sheet boundary (real line) are then :
$`𝒱^1(e;k;y)`$ $`=g_oe^\varphi e_\mu \psi ^\mu e^{ikX}(y),`$ (3.16)
$`𝒱^0(e;k;y)`$ $`=g_o(2\alpha ^{})^{1/2}e_\mu (i\dot{X}^\mu +2\alpha ^{}k\psi \psi ^\mu )e^{ikX}(y),`$ (3.17)
where the fields are all at $`y`$, and the products of fields at the same point are normal-ordered. The two operators are in -1 and 0 “pictures” - it turns out that when calculating an amplitude the sum of pictures has to be -2. Note that in addition to $`X^\mu `$ and $`\psi ^\mu `$ fields we have a $`\varphi `$ ghost field. Furthermore, we will have the $`𝒱`$’s multiplied by another ghost field $`c`$, in case their positions $`y`$ on the world sheet are *gauge-fixed*. This gauge-fixing comes about as follows. In action (3.3) the positions of the vertices are a property of the world sheet $`𝒲`$, therefore are integrated over, as different world sheets. However some of the different configurations are related by conformal transformations and are equivalent. It turns out that by fixing the gauge we can fix the positions of 3 vertices - these get multiplied by a ghost field $`c`$. The ones that are left over (one, in the case of 4 strings) have to have their positions integrated as part of the moduli $``$ of the topology in (3.4). In addition, different cyclic ordering of the three fixed coordinates on the boundary of the disk are not actually conformally equivalent, so we have to sum over the two possibilities also as part of the moduli (discreet, this time). Actually, the “pictures” and the $`\varphi `$ field come from a similar gauge fixing, but of an *anticommuting* coordinate , that we didn’t discuss.
Finally, note the undetermined constant $`g_o`$ appearing in $`𝒱`$’s, which should make the states correctly normalized. It can be related later to other unknown constants, but in the end we are still left with one constant, that quantifies the strength of open string interactions, and we choose it to be $`g_o`$. Also we have another fundamental constant of string theory $`\alpha ^{}`$.
With all these considerations the tree-level amplitude for 3 open strings looks like:
$`S_o^3`$ $`={\displaystyle 𝒟X𝒟\psi 𝒟\varphi 𝒟c\mathrm{exp}(iS[𝒲,X,\psi ,\varphi ,c])c𝒱_1^1(y_1)c𝒱_2^1(y_2)c𝒱_3^0(y_3)}`$ (3.18)
$`+(𝒱_1𝒱_2).`$
The additional term with $`𝒱_1`$ and $`𝒱_2`$ interchanged places the same strings only in cyclically reversed order on the boundary of the world sheet, which gives the other, conformally-inequivalent term. The values for $`y_i`$ can be chosen arbitrarily and the $`S_o^3`$ shouldn’t depend on that.
We note some facts about the action $`S`$ to proceed further. We see that it depends on all the fields and on the world sheet itself. The explicit $`𝒲`$-dependence of $`S`$ is just a constant factor having to do with the curvature of the world sheet. For our cases of interest the value is $`e^\lambda `$ for a disk and $`e^{2\lambda }`$ for a sphere with some constant $`\lambda `$. Furthermore, the remaining action splits into the sum of actions for different fields:
$$S_𝒲[X,\psi ,\varphi ,c]=\left(S_X[X]+S_\psi [\psi ]+S_\varphi [\varphi ]+S_c[c]\right)_𝒲,$$
(3.19)
so the integral giving the amplitude factors into integrals over different fields. We note by subscript $`𝒲`$ that the field-action is still implicitly world-sheet dependent (it matters where the fields are defined), so when we cite some result, we indicate for what world-sheet it is valid.
To shorten the notation we will write the integrals as *expectation values*. By the expectation value $`𝒪_𝒲`$ we mean only to include the integral and the action of the fields that appear in $`𝒪`$, that is:
$$𝒪_𝒲=𝒟[\mathrm{\Phi }]𝒪\mathrm{exp}(iS_\mathrm{\Phi }),$$
(3.20)
if $`\mathrm{\Phi }`$ is the set of fields that $`𝒪`$ depends on. The subscript $`𝒲`$ for the expectation value indicates the world-sheet that the fields live on.
Factoring out the explicit topology-dependent factor, we then write the amplitude (3.18) as:
$$S_o^3=e^\lambda c𝒱_1^1(y_1)c𝒱_2^1(y_2)c𝒱_3^0(y_3)_D+(𝒱_1𝒱_2),$$
(3.21)
where $`D`$ indicates the disk world-sheet. We will drop this index when it’s obvious what world-sheet we are considering. Plugging in the vertex operator expressions (3.16), (3.17) into (3.21) and factoring some fields out we get:
$`S_o^3`$ $`=e^\lambda g_o^3(2\alpha ^{})^{1/2}c(y_1)c(y_2)c(y_3)e^\varphi (y_1)e^\varphi (y_2)`$
$`\times e_{1\mu }e_{2\nu }e_{3\rho }\underset{E_X}{\underset{}{\psi ^\mu e^{ik_1X}(y_1)\psi ^\nu e^{ik_2X}(y_2)(i\dot{X}^\rho +2\alpha ^{}k_3\psi \psi ^\rho )e^{ik_3X}(y_3)}}`$
$`+(𝒱_1𝒱_2).`$ (3.22)
Here interchanging the $`𝒱`$’s just means the interchanging $`k`$’s and $`e`$’s. We can further factor $`E_X`$ which still contains two fields, $`X`$ and $`\psi `$:
$`E_X^{\mu \nu \rho }`$ $`=i\psi ^\mu (y_1)\psi ^\nu (y_2)e^{ik_1X}(y_1)e^{ik_2X}(y_2)\dot{X}^\rho e^{ik_3X}(y_3)`$
$`+2\alpha ^{}\psi ^\mu (y_1)\psi ^\nu (y_2)k_3\psi \psi ^\rho (y_3)e^{ik_1X}(y_1)e^{ik_2X}(y_2)e^{ik_3X}(y_3).`$ (3.23)
Now in order to evaluate $`S_o^3`$ we need the expectation values for these different fields on a disk. For fields $`X`$ and $`\psi `$ we will use Wick’s theorem to sum over all contractions, so we will list rules for contracting by writing
XY delimited-⟨⟩
XY \left<\mathop{\vtop{\halign{#\cr$\hfil\displaystyle{XY}\hfil$\crcr\kern 3.0pt\nointerlineskip\cr$\vrule height=3.99994pt,width=0.5pt\leaders\vrule height=0.5pt,depth=0.0pt\hfill\leaders\vrule height=0.5pt,depth=0.0pt\hfill\vrule height=3.99994pt,width=0.5pt$
\crcr\kern 3.0pt\cr}}}\limits\dots\right> to mean the term when contracting $`X`$ and $`Y`$. Also note that fields at the same point on the world-sheet are always normal-ordered, and so are not contracted with each other (the term gives zero). Then the values that we need are :
$`(\text{notation:}y_{ij}`$ $`y_iy_j)`$ (3.24)
$`{\displaystyle \underset{i}{}}e^{ik_iX(y_i)}_D`$ $`=iC_D^X(2\pi )^D\delta ^D(\mathrm{\Sigma }_ik_i){\displaystyle \underset{i<j}{}}|y_{ij}|^{2\alpha ^{}k_ik_j},`$ (3.25)
$`\stackrel{ikX(y_2)}{\begin{array}{c}\dot{X}^\mu (y_1)e\hfill \\ \text{ }\text{ }\hfill \end{array}}\mathrm{}_D`$ $`=2i\alpha ^{}{\displaystyle \frac{k^\mu }{y_{12}}}e^{ikX(y_2)}\mathrm{}_D,`$ (3.28)
$`\stackrel{\nu }{\begin{array}{c}\dot{X}^\mu (y_1)\dot{X}\hfill \\ \text{ }\text{ }\hfill \end{array}}(y_2)\mathrm{}_D`$ $`=2\alpha ^{}{\displaystyle \frac{\eta ^{\mu \nu }}{y_{12}^2}}\mathrm{}_D,`$ (3.31)
$`\stackrel{\nu }{\begin{array}{c}\psi ^\mu (y_1)\psi \hfill \\ \text{ }\text{ }\hfill \end{array}}(y_2)\mathrm{}_D`$ $`={\displaystyle \frac{\eta ^{\mu \nu }}{y_{12}}}\mathrm{}_D,\text{with }1_{\psi ,D}=1,`$ (3.34)
$`c(y_1)c(y_2)c(y_3)_D`$ $`=C_D^g|y_{12}y_{13}y_{23}|,`$ (3.35)
$`e^{\varphi (y_1)}e^{\varphi (y_2)}_D`$ $`={\displaystyle \frac{1}{|y_{12}|}},`$ (3.36)
and there exists the relation between the constants :
$$C_De^\lambda C_D^XC_D^g=\frac{1}{\alpha ^{}g_o^2},$$
(3.37)
that will allow us to leave only $`g_o`$ and $`\alpha ^{}`$ in the amplitudes. With those expectation values we can pretty straightforwardly evaluate (3.22), summing over possible contractions when needed. For example in $`E_X`$ we have:
$`\psi ^\mu (y_1)\psi ^\nu (y_2)k_3\psi \psi ^\rho (y_3)`$ $`=\underset{3}{\overset{\rho }{\begin{array}{c}\psi ^\mu _1\psi \hfill \\ \text{ }\text{ }\hfill \end{array}}}\underset{3}{\begin{array}{c}\psi ^\nu _2k_3\psi \hfill \\ \text{ }\text{ }\hfill \end{array}}\underset{3}{\begin{array}{c}\psi ^\mu _1k_3\psi \hfill \\ \text{ }\text{ }\hfill \end{array}}\underset{3}{\overset{\rho }{\begin{array}{c}\psi ^\nu _2\psi \hfill \\ \text{ }\text{ }\hfill \end{array}}}`$ (3.46)
$`={\displaystyle \frac{\eta ^{\mu \rho }}{y_{13}}}{\displaystyle \frac{k_3^\nu }{y_{23}}}{\displaystyle \frac{k_3^\mu }{y_{13}}}{\displaystyle \frac{\eta ^{\nu \rho }}{y_{23}}},`$ (3.47)
where in the first equality the minus sign comes from permuting anticommuting fields and we indicate the coordinate by subscript (we will use this notation for brevity when not ambiguous). An example of how $`X`$ fields contract in $`E_X`$ is:
$`e^{ik_1X_1}e^{ik_2X_2}\dot{X}_3^\rho e^{ik_3X_3}`$ $`=\underset{3}{\overset{\rho }{\begin{array}{c}e^{ik_1X_1}e^{ik_2X_2}\dot{X}\hfill \\ \text{ }\text{ }\hfill \end{array}}}e^{ik_3X_3}+e^{ik_1X_1}\underset{3}{\overset{\rho }{\begin{array}{c}e^{ik_2X_2}\dot{X}\hfill \\ \text{ }\text{ }\hfill \end{array}}}e^{ik_3X_3}`$ (3.52)
$`=2i\alpha ^{}\left({\displaystyle \frac{k_1^\rho }{y_{31}}}+{\displaystyle \frac{k_2^\rho }{y_{32}}}\right)e^{ik_1X_1}e^{ik_2X_2}e^{ik_3X_3}`$ (3.53)
$`=2C_D^X\alpha ^{}(2\pi )^D\delta ^D(\mathrm{\Sigma }_ik_i)\left({\displaystyle \frac{k_1^\rho }{y_{31}}}+{\displaystyle \frac{k_2^\rho }{y_{32}}}\right){\displaystyle \underset{i<j}{}}|y_{ij}|^{2\alpha ^{}k_ik_j}.`$ (3.54)
Note an important identity for 3 massless string scattering:
$$k_ik_j=0=\frac{1}{2}\left((k_i+k_j)^2k_i^2k_j^2\right),$$
(3.55)
since the sum of two momenta is minus the third one because of the delta function (momentum conservation), and $`k_i^2=0`$ for all three. Another identity, that comes from $`e_ik_i=0`$ (no sum) and $`\mathrm{\Sigma }_ik_i=0`$ is:
$$e_1k_2=e_1k_3,$$
(3.56)
and similar for other $`e_i`$’s.
We can continue to evaluate parts of $`S_o^3`$ using the above identities for simplification. In the end the $`y_i`$’s drop out, as they should (or we can fix their values from the beginning to simplify the calculation), and the result is :
$`S_o^3`$ $`=i{\displaystyle \frac{g_o}{\sqrt{2\alpha ^{}}}}(2\pi )^D\delta ^D(\mathrm{\Sigma }_ik_i)e_{1\mu }e_{2\nu }e_{3\rho }V_{123}^{\mu \nu \rho }+(𝒱_1𝒱_2),`$ (3.57)
$`V_{123}^{\mu \nu \rho }`$ $`2(\eta ^{\mu \nu }k_1^\rho +\eta ^{\nu \rho }k_2^\mu +\eta ^{\rho \mu }k_3^\nu ).`$ (3.58)
It actually turns out that the interchanged term exactly cancels the first one so $`S_o^3=0`$ (which is reasonable for the coupling of three photons). However, for closed string calculation the result without the interchange will be useful. For that we rewrite (3.57) as:
$`S_o^3`$ $`=(2\pi )^D\delta ^D(\mathrm{\Sigma }_ik_i)[A_o^3(1;2;3)+A_o^3(2;1;3)],`$ (3.59)
$`A_o^3(a;b;c)`$ $`=i{\displaystyle \frac{g_o}{\sqrt{2\alpha ^{}}}}e_{a\mu }e_{b\nu }e_{c\rho }V_{abc}^{\mu \nu \rho },`$ (3.60)
and later we will use $`A_o^3`$. Note that the $`A`$-amplitude, which we get from the S-matrix $`S`$ by factoring out $`2\pi `$’s and $`\delta `$, is exactly the analog of $`i`$ in field theory (2.43). Since in field theory the vertex amplitude is the $`i`$ for the process described by that vertex, we can directly compare amplitude $`A^3`$ from string theory with the corresponding field theory vertex. We will be able to do that for 3-closed string scattering.
Now lets move to 4-open strings. All the discussion above applies directly here, so now it is just a direct calculation. In terms of vertex operators it is:
$$S_o^4=e^\lambda dy_4c𝒱_1^1(y_1)c𝒱_2^1(y_2)c𝒱_3^0(y_3)𝒱_4^0(y_4)+(𝒱_1𝒱_2),$$
(3.61)
note that one position $`y_4`$ is not fixed, and there is no corresponding $`c`$-field. We can straightforwardly expand the expression and evaluate the expectation values. For the evaluation of the integral it is convenient to fix values $`y_1=0`$, $`y_2=1`$, $`y_3\mathrm{}`$. As for QFT 4-massless scattering, the momenta here satisfy $`k_i^2=0`$ and $`\mathrm{\Sigma }_ik_i=0`$ (we get the $`\delta `$-function again from expectation of exponentials), therefore, it is useful again to use Mandelstam variables (2.46). In the end we get the result :
$`xS_o^4=`$ $`8ig_o^2\alpha ^{}(2\pi )^D\delta ^D(\mathrm{\Sigma }_ik_i)K_oF_o,`$ (3.62)
$`K_o`$ $`{\displaystyle \frac{1}{4}}(ute_1e_2e_3e_4+(2\text{ perm.}))`$
$`{\displaystyle \frac{1}{2}}(se_3k_2e_1k_4e_2e_4+(11\text{ perm.}))e_{1\mu }e_{2\nu }e_{3\rho }e_{4\sigma }K^{\mu \nu \rho \sigma },`$ (3.63)
$`F_o`$ $`{\displaystyle \frac{\mathrm{\Gamma }(\alpha ^{}s)\mathrm{\Gamma }(\alpha ^{}t)}{\mathrm{\Gamma }(1+\alpha ^{}u)}}+{\displaystyle \frac{\mathrm{\Gamma }(\alpha ^{}t)\mathrm{\Gamma }(\alpha ^{}u)}{\mathrm{\Gamma }(1+\alpha ^{}s)}}+{\displaystyle \frac{\mathrm{\Gamma }(\alpha ^{}u)\mathrm{\Gamma }(\alpha ^{}s)}{\mathrm{\Gamma }(1+\alpha ^{}t)}}.`$ (3.64)
The permutations in the expression for $`K_o`$ mean that we should add all *inequivalent* terms that we can get from the first one by permuting indices $`\{1,2,3,4\}`$ on $`k_i`$ and $`e_i`$ simultaneously. For the first term the 2 others we get by $`(23)`$ and $`(24)`$. The second term is only identical under $`((13),(24))`$ interchange, so we have a total of $`24/2=12`$ terms, where 24 is the size of 4-permutation group. The $`\mathrm{\Gamma }`$’s in $`F_o`$ appear from $`y_4`$ integration, and they contain all the relevant poles associated with the intermediate string - we will examine such factor in more detail for the closed string. Finally, note that $`K^{\mu \nu \rho \sigma }`$ is implicitly defined in (3.63) as the factor multiplying the corresponding components of $`e_i`$’s. This definition will be used for the closed string.
### 3.4 Closed string tree amplitudes
We proceed now to our final goal - calculation of 4 dilaton scattering. We will have to do it in a couple of steps. First, we will calculate the general amplitude for 4 massless closed string scattering by another closed string. Then, by choosing the dilaton polarizations for the 4 strings we will have 4 dilaton scattering amplitude by *arbitrary* closed string. This will include more than just graviton. We will then have to analyze the amplitude in some detail to pick out the part due to the graviton.
The amplitude calculation for closed string is analogous to the open string in previous section. One difference is that this time the world-sheet is a sphere, as in Fig. 6, which will introduce some changes in expectation values. For explicit expressions we fix the sphere to be represented by the full complex plane with coordinate $`z`$. First, we need the vertex operators that for the massless closed string states (3.6) are :
$`𝒱^{1,1}(z)`$ $`=g_ce^{\varphi \stackrel{~}{\varphi }}e_{\mu \nu }\psi ^\mu \stackrel{~}{\psi }^\nu e^{ikX}(z),`$ (3.65)
$`𝒱^{0,0}(z)`$ $`={\displaystyle \frac{2g_c}{\alpha ^{}}}e_{\mu \nu }(i_zX^\mu +{\displaystyle \frac{\alpha ^{}}{2}}k\psi \psi ^\mu )(i_{\overline{z}}X^\nu +{\displaystyle \frac{\alpha ^{}}{2}}k\stackrel{~}{\psi }\stackrel{~}{\psi }^\nu )e^{ikX}(z),`$ (3.66)
where (-1,-1) and (0,0) again are different pictures and we want the total sum in the amplitude to be (-2,-2). Note the different constant of normalization $`g_c`$ which will in the end quantify the closed string coupling strength. As with the open strings, 3 coordinates of vertex operators will be gauge-fixed in an amplitude and those operators will be multiplied by ghost fields $`c\stackrel{~}{c}`$. The other coordinates will be integrated over the complex plane. Note that unlike for the disk, where there were two inequivalent cyclic orderings of fixed coordinates, on a sphere any positioning of 3 points is equivalent so there will be only one term.
We can see already that the operators are analogous to the open strings, except that there are two sets of fields $`\{\psi ,\varphi ,c\}`$ and $`\{\stackrel{~}{\psi },\stackrel{~}{\varphi },\stackrel{~}{c}\}`$, which are called left-moving and right-moving respectively. Their actions are independent so the expectation values factor. There is only one field $`X`$, but the derivatives $`_zX^\mu `$ and $`_{\overline{z}}X^\mu `$ are again independent in a sense that the expectation value between them is zero, so the expectation values for $`X`$ will also factor into left- and right-moving parts. It is true, though, that there is only one set of $`e^{ikX}`$, that will give an overall factor to the amplitude.
We won’t list the expectation values for these fields on a sphere, but they are essentially the same as (3.25)-(3.36) for the left- and right-moving parts separately, with $`y`$ replaced by $`z`$ for left- and with $`\overline{z}`$ for right-moving, and also with $`\alpha ^{}\alpha ^{}/4`$. The corresponding constants are also replaced by new ones: $`C_S^X`$ and $`C_S^g`$ with a relation :
$$C_Se^{2\lambda }C_S^XC_S^g=\frac{8\pi }{\alpha ^{}g_c^2}.$$
(3.67)
With all this in mind it is clear that the expectation value for closed string vertex operators *before integration* is two copies of open expectation values: one from left-moving and one from right-moving parts . The difference is only in constants and in that there is only one common $`\delta `$-function. The integration, in case of 4 strings, can give more non-trivial factors.
Taking care of the constant factors we get the following relations between $`A`$-amplitudes for 3 strings:
$$A_c^3=i\pi \frac{\alpha ^{}}{2}\frac{g_c}{g_o^2}\left[A_o^3\left(\frac{\alpha ^{}}{4}\right)\stackrel{~}{A}_o^3\left(\frac{\alpha ^{}}{4}\right)\right],$$
(3.68)
and for 4 strings :
$$A_c^4=i\pi ^2\alpha ^{}\frac{g_c^2}{g_o^4}\frac{1}{\mathrm{\Gamma }\left(\frac{\alpha ^{}}{4}t\right)\mathrm{\Gamma }\left(1+\frac{\alpha ^{}}{4}t\right)}\left[A_o^4(s,t;\frac{\alpha ^{}}{4})\stackrel{~}{A}_o^4(t,u;\frac{\alpha ^{}}{4})\right].$$
(3.69)
By $``$ product we mean that the polarizations from the two sides combine as
$$e_{1\mu }\stackrel{~}{e}_{1\nu }=e_{1\mu \nu },$$
(3.70)
because in the vertex operators (3.65) it is the overall polarization multiplying both parts. Then by $`\stackrel{~}{A}`$ is just meant that it contains $`\stackrel{~}{e}`$’s - the “right side” of polarizations. For 3 strings $`A_o^3(\alpha ^{}/4)`$ is just $`A_o^3(1;2;3)`$ in (3.60) with replaced $`\alpha ^{}\alpha ^{}/4`$. For 4 string amplitudes we are supposed to take as $`A_o^4`$ the contribution to (3.62) from one of the 6 possible cyclic ordering of strings on the boundary. Different orderings give different terms in $`F_o`$, with each term in (3.64) being the sum of two identical contributions. The argument in $`A_o^4`$ indicates which poles we should choose, so then:
$`A_o^4(s,t;{\displaystyle \frac{\alpha ^{}}{4}})`$ $`=ig_o^2\alpha ^{}K_o{\displaystyle \frac{\mathrm{\Gamma }(\frac{\alpha ^{}}{4}s)\mathrm{\Gamma }(\frac{\alpha ^{}}{4}t)}{\mathrm{\Gamma }(1+\frac{\alpha ^{}}{4}u)}},`$ (3.71)
$`\stackrel{~}{A}_o^4(t,u;{\displaystyle \frac{\alpha ^{}}{4}})`$ $`=ig_o^2\alpha ^{}\stackrel{~}{K}_o{\displaystyle \frac{\mathrm{\Gamma }(\frac{\alpha ^{}}{4}t)\mathrm{\Gamma }(\frac{\alpha ^{}}{4}u)}{\mathrm{\Gamma }(1+\frac{\alpha ^{}}{4}s)}},`$ (3.72)
while $`S_o^4`$ can be expressed as:
$$S_o^4=(2\pi )^D\delta ^D(\mathrm{\Sigma }_ik_i)\left(2A_o^4(s,t;\alpha ^{})+2A_o^4(t,u;\alpha ^{})+2A_o^4(u,s;\alpha ^{})\right).$$
(3.73)
Using now these relations between open and closed amplitudes we can write directly the amplitudes for 3 closed and 4 closed strings :
$`A_c^3`$ $`=i\pi g_cV_c,`$ (3.74)
$`V_c`$ $`e_{1\mu _1\nu _1}e_{2\mu _2\nu _2}e_{3\mu _3\nu _3}V^{\mu _1\mu _2\mu _3}V^{\mu _1\mu _2\mu _3},`$ (3.75)
$`A_c^4`$ $`=i\pi ^2g_c^2\alpha _{}^{}{}_{}{}^{3}K_cF_c,`$ (3.76)
$`K_c`$ $`=e_{1\mu _1\nu _1}e_{2\mu _2\nu _2}e_{3\mu _3\nu _3}e_{4\mu _4\nu _4}K^{\mu _1\mu _2\mu _3\mu _4}K^{\nu _1\nu _2\nu _3\nu _4},`$ (3.77)
$`F_c`$ $`={\displaystyle \frac{\mathrm{\Gamma }(\frac{\alpha ^{}}{4}s)\mathrm{\Gamma }(\frac{\alpha ^{}}{4}t)\mathrm{\Gamma }(\frac{\alpha ^{}}{4}u)}{\mathrm{\Gamma }(1+\frac{\alpha ^{}}{4}s)\mathrm{\Gamma }(1+\frac{\alpha ^{}}{4}t)\mathrm{\Gamma }(1+\frac{\alpha ^{}}{4}u)}},`$ (3.78)
with $`V^{\mu \nu \rho }`$ defined in (3.58) and $`K^{\mu \nu \rho \sigma }`$ in (3.63). The S-matrix amplitude for closed strings is related for both 3 and 4 as:
$$S_c^{3,4}=(2\pi )^D\delta ^D(\mathrm{\Sigma }_ik_i)A_c^{3,4}.$$
(3.79)
There is no sum over interchanges as for open strings (3.59), so $`A_c`$ gives the whole tree-level term corresponding to $`i`$ in QFT.
### 3.5 Dilaton amplitude
Now that we have a general expression for the 4 closed string scattering amplitude, we can get the one we are primarily interested in - 4 dilaton interaction. We simply have to plug in the polarization (3.11) in the expression (3.77) for $`K_c`$. The calculation is rather tedious, note that we have two copies of $`K^{\mu \nu \rho \sigma }`$, which is given by permutations in (3.63), contracted together by $`e_{\mu \nu }^{(\varphi )}`$’s:
$`K_c^{(\varphi )}`$ $`={\displaystyle \frac{1}{(D2)^2}}(\eta _{\mu _1\nu _1}k_{1\mu _1}\xi _{1\nu _1}\xi _{1\mu _1}k_{1\nu _1})(\eta _{\mu _2\nu _2}k_{2\mu _2}\xi _{2\nu _2}\xi _{2\mu _2}k_{2\nu _2})`$ (3.80)
$`\times (\eta _{\mu _3\nu _3}k_{3\mu _3}\xi _{3\nu _3}\xi _{3\mu _3}k_{3\nu _3})(\eta _{\mu _4\nu _4}k_{4\mu _4}\xi _{4\nu _4}\xi _{4\mu _4}k_{4\nu _4})K^{\mu _1\mu _2\mu _3\mu _4}K^{\nu _1\nu _2\nu _3\nu _4},`$
but again using Mandelstam variables and the momentum conservation the result in the end simplifies to:
$$K_c^{(\varphi )}=\frac{1}{16}(t^2u^2+u^2s^2+s^2t^2).$$
(3.81)
We can then put together the total tree-level amplitude for 4 dilaton scattering by *any* closed string:
$$A_c^{4(\varphi )}=\frac{i\pi ^2g_c^2\alpha _{}^{}{}_{}{}^{3}}{16}\frac{\mathrm{\Gamma }(\frac{\alpha ^{}}{4}s)\mathrm{\Gamma }(\frac{\alpha ^{}}{4}t)\mathrm{\Gamma }(\frac{\alpha ^{}}{4}u)}{\mathrm{\Gamma }(1+\frac{\alpha ^{}}{4}s)\mathrm{\Gamma }(1+\frac{\alpha ^{}}{4}t)\mathrm{\Gamma }(1+\frac{\alpha ^{}}{4}u)}(t^2u^2+u^2s^2+s^2t^2).$$
(3.82)
Again, we emphasize that this amplitude contains not only gravitational interaction. We can get insight into what particles mediate the interaction by looking at the poles of the amplitude - they appear when the intermediate particle is on-shell, therefore, we can deduce the masses of the particles that contribute to the interaction.
This information is contained in the function $`F_c`$ (note that it is polarization-independent so this analysis of poles is valid for any 4-closed string interaction). Function $`\mathrm{\Gamma }`$ has poles at:
$$|\mathrm{\Gamma }(x)|\mathrm{}x=n,n=0,1,2,\mathrm{},$$
(3.83)
and $`\mathrm{\Gamma }(x)`$ is never 0. So then $`F_c(s,t,u)`$ will have poles:
$$|F_c(s,t,u)|\mathrm{}s,t,u=\frac{4}{\alpha ^{}}n,n=0,1,2,\mathrm{},$$
(3.84)
where we mean that *any* of the three variables satisfies the condition. These are also the poles of the whole amplitude<sup>5</sup><sup>5</sup>5Note that the $`K_c`$ term in general does not cancel the pole, because if, let’s say, $`s=0`$ then $`t=u`$ and $`K_c^{(\varphi )}t^4`$. This is not 0 unless $`s=t=u=0`$, which is a special case that we ignore (it has to be taken as a limit of our amplitude).. As mentioned in the discussion of the QFT amplitude, $`s,t,u=m^2`$ \- the mass of the intermediate particle in each of the channels. Since $`A_c^4`$ takes into account all the 3 channels, it is natural that there are identical poles for each of them, and we can conclude that the particles mediating the interaction have masses:
$$m^2=\frac{4}{\alpha ^{}}n,n=0,1,2,\mathrm{}$$
(3.85)
This is exactly the spectrum of closed strings in the theory . Appart being a check that the result is consistent, this tells us that the massless closed strings couple to closed strings of any mass.
Now we want to limit the intermediate particle to being a graviton. The first thing we can do is limiting it to be *massless* \- that could be in general a graviton, an antisymmetric ensor, or another dilaton. This is done easily by taking the low-energy limit of the string theory by letting $`\alpha ^{}0`$. Since the masses of strings are proportional to $`1/\alpha ^{}`$, as in (3.85) and similarly for open strings, this limit only allows massless strings to be created, which is what we want. Expanding $`F_c`$ in series of $`\alpha ^{}`$ gives :
$$F_c=\frac{64}{\alpha _{}^{}{}_{}{}^{3}stu}+O(\alpha _{}^{}{}_{}{}^{0}).$$
(3.86)
Plugging this back in (3.82) gives:
$$A_c^{4(\varphi )}=4i\pi ^2g_c^2\left(\frac{tu}{s}+\frac{us}{t}+\frac{st}{u}\right)+O(\alpha _{}^{}{}_{}{}^{3}),$$
(3.87)
so for $`\alpha ^{}0`$ we simply get this first, $`\alpha ^{}`$-independent term. This tells us the amplitude for dilaton scattering by a *massless* closed string. We can already see the correspondence with QFT amplitude (2.54), and, actually, this $`A_c^{4(\varphi )}`$ *will* turn out to be just the gravitational amplitude, but we still have to show that.
What we want is to analyze the coupling of two dilatons with a third massless closed string, which we can do by looking at the 3-string amplitude (3.74) with appropriate polarizations. Such process is a part of 4-string amplitude in any of the channels (consider the Feynman diagrams from QFT). If we find that the scattering amplitude of two dilatons into a third particle of some kind is 0, it means that this intermediate particle can not contribute to $`A_c^{4(\varphi )}`$.
We calculate then the amplitude $`A_c^3`$ with two of the polarizations taken to be dilatons and the third one - arbitrary $`e_{\mu \nu }`$. Plugging the values into $`V_c`$ in (3.75) we have:
$`V_c^{(\varphi )}`$ $`={\displaystyle \frac{4}{D2}}e_{\mu _1\nu _1}(\eta _{\mu _2\nu _2}k_{2\mu _2}\xi _{2\nu _2}\xi _{2\mu _2}k_{2\nu _2})(\eta _{\mu _3\nu _3}k_{3\mu _3}\xi _{3\nu _3}\xi _{3\mu _3}k_{3\nu _3})`$ (3.88)
$`\times (\eta ^{\mu _1\mu _2}k_1^{\mu _3}+\eta ^{\mu _2\mu _3}k_2^{\mu _1}+\eta ^{\mu _3\mu _1}k_3^{\mu _2})(\eta ^{\nu _1\nu _2}k_1^{\nu _3}+\eta ^{\nu _2\nu _3}k_2^{\nu _1}+\eta ^{\nu _3\nu _1}k_3^{\nu _2}).`$ (3.89)
Using 3-massless particle relations $`k_ik_j=e_{i\mu \nu }k_i^\mu =e_{i\mu \nu }k_i^\nu =0`$ and momentum conservation this evaluates to:
$$V_c^{(\varphi )}=2e_{\mu \nu }(k_2^\mu k_3^\nu +k_3^\mu k_2^\nu ).$$
(3.90)
Now we can easily see what happens with different $`e_{\mu \nu }`$’s. The tensor multiplying polarization is symmetric so for antisymmetric $`e_{\mu \nu }^{(b)}`$ the amplitude is immediatly 0. For dilaton polarization $`e_{\mu \nu }^{(\varphi )}`$ we only get dot-products of $`k_i`$’s in (3.90), so it is also zero. Therefore, two dilatons couple only to a graviton, among massless states, in which case the amplitude is:
$$A_c^{3(\varphi )}=2\pi ig_ce_{\mu \nu }^{(g)}(k_2^\mu k_3^\nu +k_3^\mu k_2^\nu ).$$
(3.91)
This result also proves that the (3.87) gives exactly the gravitational scattering amplitude, which we write as the S-matrix for our final result:
$$S_{\text{grav}}^{4(\varphi )}=i(2\pi g_c)^2(2\pi )^D\delta ^D(\mathrm{\Sigma }_ik_i)\left(\frac{tu}{s}+\frac{us}{t}+\frac{st}{u}\right).$$
(3.92)
We confirm that it is equal to $`S`$ calculated in QFT in (2.54) with identification of constants :
$$\kappa =2\pi g_c.$$
(3.93)
In addition we get a second check: if we write the QFT vertex for 2-dilaton-graviton coupling $`V_{(\varphi \varphi h)}^{\mu \nu }`$ in (2.41) as an on-shell scattering amplitude:
$$i=h_{\mu \nu }V_{(\varphi \varphi h)}^{\mu \nu }=i\kappa h_{\mu \nu }(k_1^\mu k_2^\nu +k_2^\mu k_1^\nu ),$$
(3.94)
the result again matches SST amplitude $`A_c^{3(\varphi )}`$ in (3.91) with the same identification $`\kappa =2\pi g_c`$.
## 4 Conclusions
This finishes our comparison of the quantum field theory and superstring theory amplitudes. We did get a matching result in SST for the process corresponding to gravitational scattering of 4 massless scalars in QFT, and we found a relation between the constants in the two theories. As promised in the introduction this does show that we can describe gravitation by SST, therefore, this theory, being consistent as opposed to QFT of gravity, has a chance of being the unifying theory describing all particles and interactions.
Looking more generally, however, we only demonstrated this feature for a very specific case. It is possible to argue much more generally , what kind of effective actions the string theory reduces to in low-energy limit, and what those actions correspond to in QFT point of view. In that respect, the calculation in this paper is not so much valuable for the result itself, but more as an exercise, allowing us to go through many important topics in both quantum field theory and superstring theory, and demonstrating (though, a small part of) relationships between them.
### Acknowledgments
I am grateful to Professor Lüst, who agreed to supervise me for this work and the relevant studies, and who pointed this exercise problem to me. I owe many thanks to Professor Zwiebach who first introduced me to the string theory at his course at M.I.T., and encouraged me for further studies in this field.
I learned the most of the material directly relevant to this paper from: \- general relativity; \- quantum field theory; \- string theory.
I thank Max-Planck-Institut für Physik, that supported the internship, during which the work was done.
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# Boltzmann-Gibbs Distribution of Fortune and Broken Time-Reversible Symmetry in Econodynamics
## I Introduction
Carefully and extensive analysis of real economic and financial data have revealed various exponential and power-law distributions regarding to money, income, wealth, and other economic quantities in ancient and modern social societies ms ; bp ; dy ; bm . A remarkable analogy between the economic system and a thermodynamical system has been revealed in a recent study dy . Using a detailed microdynamical model with time reversal symmetry, it was demonstrated that a Boltzmann-Gibbs type distribution exists in economic systems. Indeed, ample empirical analysis support such suggestion dy ; sy . Nevertheless, different microdynamical models lead to apparently different distributions bm ; dy . Those distributions are supported by empirical data, too. The nature of such difference may reveal the difference in corresponding economic structure. For example, this difference has been tentatively attributed to the role played by time reversal symmetry in microdynamical models dy .
Following the tradition of synthesizing interdisciplinary knowledge, such as from biology, physics, and finance farmer , in this letter we argue that irrespective of the time reversal symmetry the Boltzmann-Gibbs type distribution always exists. A broader theoretical base is thus provided. The demonstration is performed within the framework of stochastic differential equations and is based on a novel mathematical structure discovered during a recent study in gene regulatory network dynamics zhu . In the light of thermodynamical description, possible explanations for the origin of the difference in various empirical distributions are proposed.
## II Boltzmann-Gibbs distribution in finance
Stochastic differential equations or Langevin equations and their corresponding Fokker-Planck equations have been shown to be a useful modelling tool in economy bs ; heston ; ms ; bp . One of the best examples is the Black-Scholes formula in option pricing theory bs . This kind of mathematical formulation provides a direct connections between the microdynamics and the stationary state and has been used to generate various distribution laws ms ; bp ; bm ; dy ; sy .
Specifically, the stochastic differential equation may take the following form:
$$\dot{𝐪}=𝐟(𝐪)+N_I(𝐪)\xi (t),$$
(1)
where $`𝐟`$ and $`𝐪`$ are $`n`$-dimensional vectors and $`𝐟`$ a nonlinear function of the state variable $`𝐪`$. The state variable is the quantity to specify the economic system. It may be the money, the income, or other suitable indices. The noise $`\xi `$ is a standard Gaussian white noise with $`l`$ independent components: $`\xi _i=0`$, $`\xi _i(t)\xi _j(t^{})=2\delta _{ij}\delta (tt^{})`$, and $`i,j=1,2,\mathrm{},l`$. In Eq.(1) we explicitly factorize out the pure noise source and the state variable dependent part for the convenient of later description.
The specification of the noise in Eq.(1) is through the $`n\times n`$ diffusion matrix $`D(𝐪)`$ by the following matrix equation
$$N_I(𝐪)N_I^\tau (𝐪)=ϵD(𝐪),$$
(2)
where $`N_I`$ is an $`n\times l`$ matrix, $`N_I^\tau `$ is its the transpose, and $`ϵ`$ a nonnegative numerical constant to keep tract of the noise. It plays the role of temperature in thermodynamics. According to Eq.(2) the $`n\times n`$ diffusion matrix $`D`$ is both symmetric and nonnegative. For the dynamics of state variable $`𝐪`$, all what needed from the noise is the diffusion matrix $`D`$. Hence, it is not necessary to require the dimension of the noise vector $`\xi `$ be the same as that of the state vector $`𝐪`$ and to require more specific knowledge of $`n\times l`$ matrices $`\{N_I\}`$ beyond Eq.(2).
It is known that even in situations that Eq.(1) is not an exact description, and perhaps it would never be in a rigorous sense in economy, it may still serve as the first approximation for further modelling ms ; bp . Indeed, it has been empirically verified to be a rather accurate description in economy bs ; heston ; rs ; loffredo . Because the energy function or Hamiltonian has played a dominant role in equilibrium physics processes, the crucial question is whether or not a similar quantity exists in a more general setting. In the following we present an argument leading to the positive answer.
There exists several ways to deal with the stochastic equations equation in the form of Eq.(1) and (2). The most commonly used are those of Ito and Stratonovich methods vankampen ; ms ; bp . However, with those methods the connection between the existence of energy function like quantity in Eq.(1) and the stationary distribution is not clear when the time reversal symmetry is broken vankampen . The difficulty for finding such potential function can be illustrated by the fact that usually $`D^1(𝐪)𝐟(𝐪)`$ cannot be written as the gradient of a scalar function vankampen in the absence of detailed balance condition or in the broken time reversal symmetry. This will become precise as we proceed.
During a recent study of the robustness of the genetic switch in a living organism zhu , it was discovered that Eq.(1) can be transformed into the following form,
$$[A(𝐪)+C(𝐪)]\dot{𝐪}=_𝐪\varphi (𝐪)+N_{II}(𝐪)\xi (t),$$
(3)
where the noise $`\xi (t)`$ is from the same source as that in Eq.(1). Here we tentatively name $`A(𝐪)`$ the adaptation matrix, $`C(𝐪)`$ the conservation matrix, and the scalar function $`\varphi (𝐪)`$ the fortune function. The gradient operation in state space is denoted by $`_𝐪`$. The adaptation matrix $`A(𝐪)`$ is defined through the following matrix equation
$$N_{II}(𝐪)N_{II}^\tau (𝐪)=ϵA(𝐪),$$
(4)
which guarantees that $`A`$ is both symmetric and nonnegative. The $`n\times n`$ conservation matrix $`C`$ is antisymmetric. We define
$$A(𝐪)+C(𝐪)=1/[D(𝐪)+Q(𝐪)]M(𝐪).$$
with the $`n\times n`$ matrix $`M`$ is the solution of following two matrix equations ao04 . The first equation is the potential condition
$$_𝐪\times [M(𝐪)𝐟(𝐪)]=0,$$
(5)
which gives $`n(n1)/2`$ conditions \[the wedge product for two arbitrary vectors $`𝐯_1`$ and $`𝐯_2`$ in n-dimension: $`[𝐯_1\times 𝐯_2]_{ij}=v_{1}^{}{}_{i}{}^{}v_{2}^{}{}_{j}{}^{}v_{1}^{}{}_{j}{}^{}v_{2}^{}{}_{i}{}^{},i,j=1,2,\mathrm{},n`$ \]. The second equation is the generalized Einstein relation between the adaptation and diffusion matrices in the presence of conservation matrix
$$M(𝐪)D(𝐪)M^\tau (𝐪)=\frac{1}{2}[M(𝐪)+M^\tau (𝐪)].$$
(6)
which gives $`n(n+1)/2`$ conditions ao04 . The fortune function $`\varphi (𝐪)`$ is connected to the deterministic force $`𝐟(𝐪)`$ by
$$_𝐪\varphi (𝐪)=M(𝐪)𝐟(𝐪).$$
For simplicity we will assume $`det(A)0`$ in the rest of the letter. Hence $`det(M)0`$ kat . Thus, the adaptation matrix $`A`$, the conservation matrix $`Q`$ and the fortune function $`\varphi `$ in Eq.(3) and (4) can be completely determined by Eq.(1) and (2). The breakdown of detailed balance condition or the time reversal symmetry is represented by the finiteness of the conservation matrix
$$C(𝐪)0,$$
(7)
or equivalently $`Q0`$. The usefulness of the formulation of Eq.(3) and (4) is already manifested in the successful solution of outstanding stable puzzle in gene regulatory dynamics zhu and in solving two fundamental controversies in population genetics ao05 .
A few remarks on Eq.(3) are in order. In the light of classical mechanics in physics, Eq.(3) is in precisely the form of Langevin equation. The fortune function $`\varphi `$ corresponds to the potential function but opposite in sign to reflect the fact that in economy there is a tendency to seek the peak or maximum of fortune. The adaptive matrix $`A`$ plays the role of friction. It represents adaptive dynamics and is the dynamical mechanism to seek the nearby fortune peak. The conservation matrix $`C`$ plays the role analogous to a magnetic field. Its dynamics is similar to that of the Lorentz force, hence conserves the fortune. As in classical mechanics, the finiteness of the conservation matrix $`C`$ breaks the time reversal symmetry.
It was heuristically argued ao04 and rigorous demonstrated yin that the stationary distribution $`\rho (𝐪)`$ in the state space is, if exists,
$$\rho (𝐪)\mathrm{exp}\left(\frac{\varphi (𝐪)}{ϵ}\right).$$
(8)
Therefore, the fortune function $`\varphi `$ acquires both the dynamical meaning through Eq.(3) and the steady state meaning through Eq.(8). Specifically, in the so-called zero-mass limit to differentiate from Ito and Stratonovich methods, the Fokker-Planck equation for the probability distribution $`\rho (𝐪,t)`$ takes the form yin
$$_t\rho (𝐪,t)=_𝐪^\tau M^1(𝐪)[ϵ_𝐪_𝐪\varphi (𝐪)]\rho (𝐪,t).$$
(9)
Here $`_t`$ is a derivative with respect to time and $`_𝐪`$ represents the gradient operation in state space. We note that Eq.(8) is a stationary solution to Eq.(9) even it may not be normalizable, that is, even when the partition function $`𝒵=d^n𝐪\rho (𝐪)`$ is ill-defined. Again, we emphasize that no time reversal symmetry is assumed in reaching this result. This completes our demonstration on the existence of the Boltzmann-Gibbs distribution in economy.
Using $`M^1(𝐪)=D(𝐪)+Q(𝐪)`$ and $`𝐟(𝐪)=[D(𝐪)+Q(𝐪)]_𝐪\varphi (𝐪)`$, Eq.(9) can be rewritten in a more suggestive form yin
$$_t\rho (𝐪,t)=_𝐪^\tau [ϵD(𝐪)_𝐪+ϵ(_𝐪^\tau Q(𝐪))𝐟(𝐪)]\rho (𝐪,t).$$
(10)
It is clear that in the presence of time reversal symmetry, i.e. $`Q=0`$, one can directly read the fortune function $`\varphi `$ from above form of Fokker-Planck equation as $`_𝐪\varphi (𝐪)=D^1(𝐪)𝐟(𝐪)`$.
For the sake of completeness, we list the Fokker-Placnk equations corresponding to Ito and Stratonovich treatments of Eq.(1) and (2) vankampen :
$$_t\rho _I(𝐪,t)=\underset{i=1}{\overset{n}{}}_{q_i}\left[f_i(𝐪)+ϵ\underset{j=1}{\overset{n}{}}_{q_j}D_{i,j}(𝐪)\right]\rho _I(𝐪,t),(Ito)$$
(11)
and
$$_t\rho _S(𝐪,t)=\underset{i=1}{\overset{n}{}}_{q_i}[f_i(𝐪)+\underset{j=1}{\overset{n}{}}\underset{k=1}{\overset{l}{}}N_{I}^{}{}_{ik}{}^{}(𝐪)_{q_j}N_{I}^{}{}_{jk}{}^{}(𝐪)]\rho _S(𝐪,t).(Str)$$
(12)
The connection of fortune function with both dynamical (Eq.(1) and (2)) and stationary state is indeed not clear in above two equations. Nevertheless, it has been shown yin that there are corresponding fortune functions, adaptation and conservation matrices to Eq.(11) and (12). We point out here that when the matrix $`N_I`$ is independent of state variable, Eq.(11) and (12) are the same but may still differ from Eq.(9), because the gradient of the antisymmetric matrix $`Q(𝐪)`$ may not be zero. This last property shows that the time reversal symmetry is indeed important.
## III Two Examples
The Fokker-Planck equation used by Silva and Yakovenko sy has the form:
$$_t\rho _{sy}(q,t)=_q[a(q)+_qb(q)]\rho _{sy}(q,t),$$
(13)
and Fokker-Planck equation used by Bouchaud and Mezard bm has the form
$$_t\rho _{bm}(q,t)=_q[(J(q1)+\sigma ^2q+\sigma ^2q_qq]\rho _{bm}(q,t).$$
(14)
They are all in one dimension. We immediately conclude that the conservation matrix $`C`$, equivalently $`Q`$, is zero, because there is no conservation matrix in one dimension. In accordance with the definition in the present letter, which is consistent with that in nonequilibrium processes vankampen , the dynamics described by above two equations can be effectively classified as time reversal symmetric.
Rewriting them in symmetric form with respect to the derivative of state variable $`q`$ as in Eq.(10), we have
$$_t\rho _{sy}(q,t)=_q[a(q)+(_qb(q))+b(q)_q]\rho _{sy}(q,t),$$
(15)
and
$$_t\rho _{bm}(q,t)=_q[(J(q1)+\sigma ^2q+\sigma ^2q+\sigma ^2q_q]\rho _{bm}(q,t).$$
(16)
The corresponding wealth functions can be immediate read out as
$`\varphi _{sy}(q)`$ $`=`$ $`{\displaystyle _{q_0}^q}𝑑q^{}{\displaystyle \frac{a(q^{})+(_q^{}b(q^{}))}{b(q^{})}}`$ (17)
$`=`$ $`{\displaystyle \frac{a}{b}}(qq_0),ifa,b=constant`$
$`\varphi _{bm}(q)`$ $`=`$ $`{\displaystyle _{q_0}^q}𝑑q^{}{\displaystyle \frac{(J(q^{}1)+2\sigma ^2q^{}}{\sigma ^2q_{}^{}{}_{}{}^{2}}}`$ (18)
$`=`$ $`{\displaystyle \frac{J}{\sigma ^2}}{\displaystyle \frac{1}{q}}\left(2+{\displaystyle \frac{J}{\sigma ^2}}\right)\mathrm{ln}q+{\displaystyle \frac{J}{\sigma ^2}}{\displaystyle \frac{1}{q_0}}+\left(2+{\displaystyle \frac{J}{\sigma ^2}}\right)\mathrm{ln}q_0`$
They are exactly what found in Ref.\[sy \] and \[bm \]: the first one corresponds to an exponential distribution and the second one a power law distribution according to the Boltzmann-Gibbs distribution Eq.(8).
## IV Ensembles and state variables
Having defined a precise meaning of time reversal symmetric and have demonstrated that the Boltzmann-Gibbs distribution even in the absence of time reversal symmetry, we explore further connection between econodynamics and statistical physics.
There are two general types of situations which would generate different distributions in statistical physics and thermodynamics. The first one is to link to constraints on the system under various conditions. In statistical physics such constraints are described by various ensembles and free energies. For example, there are canonical and grand-canonical ensembles. There are Gibbs and Helmholtz free energies, entropy, enthalpy, etc. Those ensembles have their characteristic distributions. It would be interesting to know the corresponding situations in economy.
Even with a given constraint, the form of distribution depends on the choice of state variable. For example, for ideal gas model, the distributions are different if views from the kinetic energy and from velocity. Hence, there is a question of appropriate state variable for a given situation, with which the physics becomes particular transparent. It would be interesting to know what be the appropriate variables to describe an economic system. Within this context, the difference between what discovered by Dragulecu and Yakovenko dy and Bouchaud and Mezard bm is perhaps more due to the difference in choices of state variables, because it seems they are describing the same situation of same system under same constraints.
A remark on terminology is in order. It was demonstrated above that regardless of the time reversal symmetry the Boltzmann-Gibbs distribution, Eq.(8), exists. The fortune function $`\varphi `$ has an additional dynamical meaning defined in Eq.(3). Both exponential and power law distribution can be represented by Eq.(8). In fact, it is well known that power law distributions exist in statistical physics. A nontrivial example is the Kosterlitz-Thouless transition kt . Thus it does not appear appropriate to call the power law distribution non-Boltzmann-Gibbs distribution. Such a terminology confusion was already noticed before ls .
In the view of the dominant role of entropy in Kosterlitz-Thouless transition kt , the ubiquitous existence of power law distribution in economy may suggest that the entropy effect is rather important in econodynamics. This may corroborate with the suggestion of “superthermal” in economy sy .
## V Conclusions
In this letter we demonstrate that the existence of Botlzmann-Gibbs distribution in finance is independent of time reversal symmetry. Both power law and exponential distributions are within its description. In analogous to similar situation in statistical physics, the differences among those distributions discovered empirically in economy are likely the result of different choices of state variables to describe the same system in econodynamics.
This work was supported in part by USA NIH grant under HG002894.
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# State-Selective Detection of Near-Dissociation Ultracold KRb 𝑋¹Σ⁺ and 𝑎³Σ⁺ Molecules
## I Introduction
Following in the footsteps of ultracold atoms, the area of ultracold molecules has developed rapidly in recent years Bahns et al. (2000); Masnou-Seeuws and Pillet (2001); Bethlem and Meijer (2003). Polar molecules, with their permanent electric dipole moments, have drawn particular attention Special Issue: “Ultracold Polar Molecules: Formation and Collisions” . The dipole moments of these heteronuclear molecules allow manipulation with applied electric fields. Because the dipole-dipole potential is both long-range and anisotropic, interactions between polar molecules are fundamentally different from those in homonuclear systems. Potential applications include novel quantum degenerate systems Santos et al. (2000); Yi and You (2000); Goral et al. (2000); Goral and Santos (2002); Damski et al. (2003); Goral et al. (2002); Baranov et al. (2002), quantum computation DeMille (2002), and tests of fundamental symmetries Kozlov and Labzowsky (1995). Investigations of collisions and reactions at these extremely low temperatures also promise to open new areas of ultracold chemistry Bohn (2001); Bodo and Gianturco (2002); Balakrishnan and Dalgarno (2001).
The technique of ultracold atom photoassociation (PA) Stwalley and Wang (1999) has been the primary means for producing molecules at translational temperatures below 1 mK Masnou-Seeuws and Pillet (2001). In this process, laser light resonantly binds two colliding atoms into an excited molecule which subsequently decays by spontaneous emission. This method has recently been successfully adapted to heteronuclear systems Kerman et al. (2004); Sage et al. ; Bergeman et al. (2004); Mancini et al. (2004); Shaffer et al. (1999); Haimberger et al. (2004); Wang (2003); Wang et al. (2004); D. Wang, J. Qi, M. F. Stone, O. Nikolayeva, B. Hattaway, S. D. Gensemer, H. Wang, W. T. Zemke, P. L. Gould, E. E. Eyler, W. C. Stwalley . Feshbach resonances have been used to make homonuclear molecules at even lower temperatures Herbig et al. (2003); Chin et al. (2003); Dürr et al. (2004); Xu et al. (2003); Jochim et al. (2003); Cubizolles et al. (2003); Strecker et al. (2003); Zwierlein et al. (2003); Regal et al. (2003); Greiner et al. (2003), in some cases under quantum degenerate conditions. Based on recent observations of Feshbach resonances in heteronuclear systems Stan et al. (2004); Inouye et al. (2004), it can be expected that this route will also yield ultracold heteronuclear molecules in due course. Both PA and Feshbach resonances tend to produce molecules in high vibrational levels with outer turning points at long range. Although these high-$`v^{\prime \prime }`$ states are of interest for some studies, many applications require low-$`v^{\prime \prime }`$ states, because of their improved stability against inelastic processes and their larger dipole moments. Various schemes for populating low-$`v^{\prime \prime }`$ states have been proposed Band and Julienne (1995); DeMille (2002); Bergeman et al. (2004); Kotochigova et al. (2004); Stwalley (2004); Damski et al. (2003), and some successful implementations have been reported Nikolov et al. (2000); Sage et al. .
Direct detection of ultracold molecules has generally utilized photoionization combined with time-of-flight mass spectroscopy. State selectivity is essential to many applications and for diagnosing transfer of population. For example, measuring the $`v^{\prime \prime }`$ dependence of vibrational quenching rates due to collisions with ultracold atoms would be an important first step in studying ultracold molecule collisions. To date, the ability to unambiguously identify the vibrational state of ultracold molecules has been limited, particularly for levels near dissociation. Two-photon one-color ionization spectra have been reported in Rb<sub>2</sub> Gabbanini et al. (2000); Fioretti et al. (2001); Kemmann et al. (2004), but definitive assignments and initial state identifications were not made. We have recently demonstrated state-selective detection of Rb<sub>2</sub>, details of which will be reported elsewhere Huang et al. . In Cs<sub>2</sub>, ultracold molecule detection spectra have been compared to absorption measurements, revealing the role of the diffuse bands in the detection process Dion et al. (2002). In K<sub>2</sub>, two-photon two-color ionization spectra enabled some identification of the low-$`v^{\prime \prime }`$ states produced by both one-photon Nikolov et al. (1999) and two-photon Nikolov et al. (2000) PA. Dissociation with a separate cw laser helped to clarify the assignment Nikolov et al. (2000). In Na<sub>2</sub>, high-resolution cw ionization spectroscopy yielded clear identification of the highest-lying bound and quasibound rovibrational-hyperfine states Fatemi et al. (2002).
To date, the only heteronuclear system in which state selectivity has been achieved is RbCs. Using two-photon two-color ionization, ultracold molecules were detected in specific vibrational levels of the $`a^3\mathrm{\Sigma }^+`$ state Kerman et al. (2004); Bergeman et al. (2004). In addition, these molecules have been transferred by stimulated emission pumping to the $`X^1\mathrm{\Sigma }^+`$ state ($`v^{\prime \prime }`$=0,1) and detected in these states with vibrational selectivity Sage et al. .
In the present work, we report on vibrationally state-selective detection of ultracold KRb molecules in high-$`v^{\prime \prime }`$ levels of both the ground $`X^1\mathrm{\Sigma }^+`$ state and the metastable $`a^3\mathrm{\Sigma }^+`$ state. These molecules are formed by cold-atom photoassociation followed by radiative decay, as shown in Fig. 1. Two-photon one-color ionization proceeds through resonant intermediate levels of the $`4^1\mathrm{\Sigma }^+`$, $`5^1\mathrm{\Sigma }^+`$, $`4^3\mathrm{\Sigma }^+`$, and $`3^3\mathrm{\Pi }`$ states, allowing vibrational state identification and determination of the relative populations. Using wavelengths in the range of 582-625 nm, we have observed vibrational levels $`v^{\prime \prime }`$=86-92 for the $`X`$ state and $`v^{\prime \prime }`$=17-23 for the $`a`$ state. We have also analyzed spectroscopy of the excited states, details of which will be reported elsewhere Wang et al. .
## II Experiment
Details of the experimental setup have been described previously Wang et al. (2004); D. Wang, J. Qi, M. F. Stone, O. Nikolayeva, B. Hattaway, S. D. Gensemer, H. Wang, W. T. Zemke, P. L. Gould, E. E. Eyler, W. C. Stwalley . Here we recount it briefly, focusing on the ionization detection. The KRb PA takes place in overlapping clouds of ultracold <sup>39</sup>K and <sup>85</sup>Rb. High atomic densities, estimated at $`3\times 10^{10}`$ cm<sup>-3</sup> for K and $`1\times 10^{11}`$ cm<sup>-3</sup> for Rb, are achieved using “dark-SPOT” MOTs for each species. Temperatures for K and Rb of 300 $`\mu `$K and 100 $`\mu `$K, respectively, are expected.
The PA process is driven with a cw tunable titanium-sapphire laser (Coherent 899-29). Its output, typically $`>`$400 mW, is focused into the overlapping MOT clouds. PA spectra, as described previously Wang et al. (2004); D. Wang, J. Qi, M. F. Stone, O. Nikolayeva, B. Hattaway, S. D. Gensemer, H. Wang, W. T. Zemke, P. L. Gould, E. E. Eyler, W. C. Stwalley , are obtained by scanning this laser and measuring the ionization signal from molecules which have radiatively decayed into the $`X^1\mathrm{\Sigma }^+`$ ground state and the metastable $`a^3\mathrm{\Sigma }^+`$ state.
In our earlier work Wang et al. (2004); D. Wang, J. Qi, M. F. Stone, O. Nikolayeva, B. Hattaway, S. D. Gensemer, H. Wang, W. T. Zemke, P. L. Gould, E. E. Eyler, W. C. Stwalley , the ionization detection used a pulsed laser with a broader linewidth and significant amplified spontaneous emission (ASE). This prevented state-selective detection but had the benefit of yielding at least some ion signal at most wavelengths, thereby facilitating the location of PA resonances.
In the present work, ionization detection is achieved with a pulsed dye laser (Continuum ND6000) pumped by a frequency-doubled Nd:YAG laser at a 10 Hz repetition rate. The use of two dyes, R610 and DCM, provides spectral coverage over the range 582 nm to 625 nm. The 0.05 cm<sup>-1</sup> linewidth of the dye laser is sufficient to resolve the vibrational structure, but not the rotational structure, of transitions from levels of the $`X^1\mathrm{\Sigma }^+`$ and $`a^3\mathrm{\Sigma }^+`$ states. This detection laser, with a pulse width of 7 ns and a typical output power of 3 mJ, is focused to a diameter of 1 mm. This is significantly larger than the 0.3 mm diameter of the MOT clouds in order to illuminate a larger fraction of the ballistically expanding cloud of cold molecules. Ions from the laser pulse are accelerated to a channeltron ion detector. KRb<sup>+</sup> is discriminated from other species (K<sup>+</sup>, Rb<sup>+</sup>, Rb$`{}_{}{}^{+}{}_{2}{}^{}`$) by its time of flight (TOF).
## III Detection Spectra for Singlet Molecules
Detection spectra are obtained by fixing the PA frequency on a resonance and recording the KRb<sup>+</sup> ion signal while scanning the pulsed laser. A 400 cm<sup>-1</sup> scan for $`X^1\mathrm{\Sigma }^+`$ state molecules is shown in Fig. 2. The spectra display structure on two scales. On a gross scale ($``$20 cm<sup>-1</sup>), nearly periodic spacings correspond to vibrational levels of the upper state of the detection transition, specifically the $`4^1\mathrm{\Sigma }^+`$ state for the spectrum shown in Fig. 2. The spectroscopy of this upper state will be described separately Wang et al. . Measured level spacings and the range of levels observed are both in very good agreement with calculations based on $`ab`$ $`initio`$ potentials Rousseau et al. (2000). We note that the increased complexity of the spectrum at higher frequencies is likely due to the appearance of transitions to the $`5^1\mathrm{\Sigma }^+`$ state, which are not yet assigned.
On a finer scale ($``$5 cm<sup>-1</sup>), the structure corresponds to the spacings between high-lying vibrational levels of the $`X^1\mathrm{\Sigma }^+`$ ground state. Specific $`X`$-state levels are assigned by matching the measured spacings with those calculated from the potential Zemke and Stwalley (2004). Examples are shown in Fig. 3 for PA to two different vibrational levels of the $`3(0^+)`$ state. As can be seen, the agreement between measured and predicted vibrational spacings is excellent, allowing unambiguous identification of the ground-state levels. A previously measured spacing between $`v^{\prime \prime }`$=86 and 87 Amiot and Vergès (2000) is also in excellent agreement. We note that although the $`3(0^+)`$ state comes from the K(4$`s`$) + Rb(5$`p_{3/2}`$) asymptote, the PA laser is tuned below the K(4$`s`$) + Rb(5$`p_{1/2}`$) asymptote, so predissociation does not occur.
The relative populations of the various $`X`$-state vibrational levels are determined by the Franck-Condon factors (FCFs) for decay of the level formed by photoassociation. We expect that as PA occurs to more deeply bound levels, with smaller outer turning points, the FCFs will favor decay to more deeply bound levels of the $`X`$-state. This is indeed the case, as seen in comparing Figs. 3(a) and 3(b). The peak of the $`X`$-state distribution shifts to a lower $`v^{\prime \prime }`$ for the larger (more negative) PA detuning. If we assume that the first (resonant) step of the two-photon detection process is saturated, and that the second (ionization) step is saturated and/or structureless, then the detection spectra will be a direct measure of the relative population of $`X`$-state vibrational levels. This distribution is given by the FCFs for decay of the photoassociated level, in this case a level of the $`3(0^+)`$ state. These FCFs, calculated using the LEVEL program LeRoy , are superimposed on the spectra in Figs. 3(a) and 3(b). Not only does the measured peak of the distribution shift with detuning as predicted, but the relative peak heights within each distribution actually match the FCF calculations rather well.
Using PA detunings for the $`3(0^+)`$ state over the range 246 cm<sup>-1</sup> to 320 cm<sup>-1</sup> (measured relative to its K(4$`S`$) + Rb(5$`P_{3/2}`$) asymptote, which lies 237.60 cm<sup>-1</sup> above the K(4$`S`$) + Rb(5$`P_{1/2}`$) asymptote), we observe $`X`$-state levels from $`v^{\prime \prime }`$=86-92, which have binding energies from 29.76 cm<sup>-1</sup> to 4.45 cm<sup>-1</sup>, respectively. The highest level in the $`X`$ state is predicted Zemke and Stwalley (2004) to be $`v^{\prime \prime }`$=98. The molecules we detect are produced over a time interval of several milliseconds before the detection laser pulse. During this time, they are exposed to the Rb MOT trap (and repump) light, tuned near the K(4$`S`$) + Rb(5$`P_{3/2}`$) asymptote, and PA light. This could cause off-resonant reexcitation and subsequent dissociation of the $`X`$-state molecules, particularly for high $`v^{\prime \prime }`$. Such a state-dependent destruction mechanism would modify the $`X`$-state vibrational distribution from that intially produced according to the FCFs. However, we see no evidence for this alteration.
## IV Detection Spectra for Triplet Molecules
State-selective detection of triplet $`a^3\mathrm{\Sigma }^+`$ molecules is carried out in a similar manner to that described above for singlet $`X^1\mathrm{\Sigma }^+`$ molecules. The only difference is that the PA laser is tuned to a level which decays to the $`a`$ state. For this work, we use the 3($`0^{}`$) state from the K(4$`S`$) + Rb(5$`P_{3/2}`$) asymptote, which correlates to the $`1^3\mathrm{\Pi }`$ state at short range D. Wang, J. Qi, M. F. Stone, O. Nikolayeva, B. Hattaway, S. D. Gensemer, H. Wang, W. T. Zemke, P. L. Gould, E. E. Eyler, W. C. Stwalley . Higher detection laser intensities are required for comparable signal levels from triplet molecules. The triplet detection spectra are obtained over the same wavelength range, but of course involve different upper states: $`4^3\mathrm{\Sigma }^+`$ and probably $`3^3\mathrm{\Pi }`$. The spectroscopy of the $`4^3\mathrm{\Sigma }^+`$ state, along with that of the $`4^1\mathrm{\Sigma }^+`$ state (used for singlet detection), will be described in a separate publication Wang et al. . A $``$400 cm<sup>-1</sup> scan is shown in Fig. 4. The structure repeating at $``$50 cm<sup>-1</sup> corresponds to the vibrational spacing of the $`4^3\mathrm{\Sigma }^+`$ state. Assigned vibrational levels are indicated in the figure. The congested region of larger signal to the blue may involve the $`3^3\mathrm{\Pi }`$ state, but specific assignments have not yet been made.
Two-photon transitions to atomic Rb Rydberg states sometimes can be seen as well in the triplet spectra, due to leakage from the space-charge-broadened Rb<sup>+</sup> TOF peak into the KRb<sup>+</sup> TOF peak. We also see ions at the Rb atomic 5p$`{}_{3/2}{}^{}`$5f one-photon transition. The 5p<sub>3/2</sub> level is populated by the MOT lasers and this dipole-forbidden transition is enabled by the $``$160 V/cm electric field used for extraction. These atomic lines serve as useful frequency markers and verify that the laser frequency measurements are accurate to 0.16 cm<sup>-1</sup>. These lines are not readily observable in the spectra of singlet molecules (Fig. 2) because lower laser intensities are used.
Fig. 5 is an expanded view of the spectrum, encompassing only one vibrational level of the $`4^3\mathrm{\Sigma }^+`$ upper state. The structure here is due to the near-dissociation levels of the $`a^3\mathrm{\Sigma }^+`$ state. The vibrational levels show spacings of 2.4 to 5.6 cm<sup>-1</sup>. Such spacings correspond to levels $`v^{\prime \prime }`$=20 to 26 in the $`ab`$ $`initio`$ potential of Kotochigova et al. Kotochigova et al. (2003) and $`v^{\prime \prime }`$=18 to 24 in the $`ab`$ $`initio`$ potentials of Park et al. Park et al. (2000) and of Rousseau et al. Rousseau et al. (2000). These three potentials have been compared by Zemke et al. Zemke et al. (2005). However, a definitive assignment has very recently become available, based on new Fourier Transform Spectra in Hannover Tiemann and Docenko . This clearly indicates we have observed levels $`v^{\prime \prime }`$=17 to 23, as shown in Table 1. We also plan to calculate improved Franck-Condon factors for the Hannover $`a^3\mathrm{\Sigma }^+`$ potential once it is available. The FCFs in Fig. 5 are based on the Kotochigova et al. potential Kotochigova et al. (2003) for levels 20 to 26. Finally, we plan direct measurement of the binding energies by scanning a separate cw laser to deplete the ground-state levels. We have recently observed this “ion dip” spectroscopy for $`X^1\mathrm{\Sigma }^+`$ $`v^{\prime \prime }`$=89.
Table I. Level spacings ($`\mathrm{\Delta }G_{v+1/2}`$, in cm<sup>-1</sup>) in the <sup>39</sup>K<sup>85</sup>Rb $`a^3\mathrm{\Sigma }^+`$ state
$`v`$ Fourier Transform Spectra PA-REMPI(This work) 17 5.49 5.60 18 4.84 4.90 19 4.23 4.23 20 3.60 3.56 21 2.99 2.96 22 2.41 2.38
As for the singlet molecules, we can use the peak heights as a measure of the lower level ($`a^3\mathrm{\Sigma }^+`$) populations. There is a small, but noticeable shift in the distribution to lower $`v^{\prime \prime }`$ for larger (more negative) PA detunings. Also shown in these figures are the calculated FCFs for decay from each 3($`0^{}`$) PA level to various levels of the $`a^3\mathrm{\Sigma }^+`$ state. Calculations using the LEVEL program LeRoy reproduce the overall locations of the distributions, including the shift with PA detuning. However, compared to the singlet spectra (Fig. 3), individual peak heights are not as accurately predicted. We hope future calculations based on the Hannover $`a^3\mathrm{\Sigma }^+`$ potential, once it is available, can give more accurate results. Linewidths in the triplet spectra are somewhat broader than those in the singlet spectra. Power broadening, spin-spin, second-order spin-orbit, rotational and hyperfine structure should all contribute to the line shapes observed.
PA detunings for the 3($`0^{}`$) state from $`244`$ cm<sup>-1</sup> to 323 cm<sup>-1</sup> (measured relative to its K(4$`P`$) + Rb(5$`P_{3/2}`$) asymptote) have been used to observe $`a`$-state levels from $`v^{\prime \prime }`$=17-23. These have binding energies from 29.02 cm<sup>-1</sup> to 5.31 cm<sup>-1</sup>, respectively. This numbering, based on the Fourier Transform Spectra in Hannover Tiemann and Docenko , is definitive because vibrational assignment is based on two different isotopes. None of the three sets of $`ab`$ $`initio`$ potentials from Kotochigova et al. (2003); Rousseau et al. (2000); Park et al. (2000) can give this exact numbering, although a good vibrational spacing match can be found if we adjust their numbering by one or three Zemke et al. (2005).
An important difference between singlet and triplet molecules is that the triplets have a non-zero magnetic moment and can therefore be magnetically trapped. We have previously demonstrated this trapping in the quadrupole magnetic field of the MOT Wang et al. (2004); D. Wang, J. Qi, M. F. Stone, O. Nikolayeva, B. Hattaway, S. D. Gensemer, H. Wang, W. T. Zemke, P. L. Gould, E. E. Eyler, W. C. Stwalley by delaying the molecule detection with respect to the turn-off of the PA laser. This difference in magnetic properties could be utilized as a singlet/triplet “filter$`\mathrm{"}`$ to distinguish the two types of molecules. We do see cases where both $`X^1\mathrm{\Sigma }^+`$ and $`a^3\mathrm{\Sigma }^+`$ molecules appear in the same region of the detection spectrum. An example, using PA to the 3($`0^+`$) state, is shown in Fig. 6. At long range, the 3($`0^+`$) state should decay to both $`X^1\mathrm{\Sigma }^+`$ and $`a^3\mathrm{\Sigma }^+`$ states D. Wang, J. Qi, M. F. Stone, O. Nikolayeva, B. Hattaway, S. D. Gensemer, H. Wang, W. T. Zemke, P. L. Gould, E. E. Eyler, W. C. Stwalley , so we expect to detect both. However, the FCFs for the first step of the detection process play an important role. For high-$`v^{\prime \prime }`$ levels of $`X^1\mathrm{\Sigma }^+`$, the overlap with $`4^1\mathrm{\Sigma }^+`$ levels comes primarily from the outer turning point. On the other hand, overlap of high-$`v^{\prime \prime }`$ levels of $`a^3\mathrm{\Sigma }^+`$ with $`4^3\mathrm{\Sigma }^+`$ levels comes primarily from the inner turning point and becomes more favorable in the lower energy region of the spectrum. The overall detection efficiencies (including the ionization step) for $`X^1\mathrm{\Sigma }^+`$ and $`a^3\mathrm{\Sigma }^+`$ become comparable in the region shown in Fig. 6. Although we do see triplet features in a primarily singlet detection spectrum (using PA to 3($`0^+`$)), we do not see singlet features in the triplet detection spectra (using PA to 3($`0^{}`$)). This is consistent with the fact that at long range, 3($`0^+`$) can decay to both to both $`X^1\mathrm{\Sigma }^+`$ and $`a^3\mathrm{\Sigma }^+`$ states, while 3($`0^{}`$) can decay only to $`a^3\mathrm{\Sigma }^+`$ D. Wang, J. Qi, M. F. Stone, O. Nikolayeva, B. Hattaway, S. D. Gensemer, H. Wang, W. T. Zemke, P. L. Gould, E. E. Eyler, W. C. Stwalley .
## V Conclusions
In summary, we have realized vibrationally state-selective detection of near-dissociation levels of ultracold KRb molecules in the $`X^1\mathrm{\Sigma }^+`$ ground state and the $`a^3\mathrm{\Sigma }^+`$ metastable state. This state-selectivity will be crucial to future experiments in ultracold molecular collisions and reactions where specific initial and final states must be measured. This capability is equally important in diagnosing population transfer, e.g., from high-$`v^{\prime \prime }`$ to low-$`v^{\prime \prime }`$ Sage et al. . In fact, the first step in our singlet detection ($`4^1\mathrm{\Sigma }^+X^1\mathrm{\Sigma }^+`$) is a good candidate for realizing this type of transfer. As an example, if we start in $`X^1\mathrm{\Sigma }^+`$ ($`v^{\prime \prime }`$=89), the FCF for excitation to $`4^1\mathrm{\Sigma }^+`$ ($`v^{}`$=40) is quite large (0.02) due to overlap at the outer turning points. On the other hand, overlap at the inner turning points gives a favorable FCF of $``$0.01 for decay (or stimulated emission) of $`4^1\mathrm{\Sigma }^+`$ ($`v^{}`$=40) to the absolute ground state, $`X^1\mathrm{\Sigma }^+`$ ($`v^{\prime \prime }`$=0). These large FCFs indicate that coherent two-photon transfer, such as STIRAP Bergmann et al. (1998), should be feasible with narrow-linewidth quasi-cw lasers.
The two-photon, one-color detection we have employed is particularly convenient because only one tunable laser is required. However, the second (ionizing) step generally requires high intensity, resulting in power broadening of the first (bound-bound) step. Two-photon, two-color detection (e.g., through states from the K(4$`s`$)+Rb(5$`p`$) asymptotes) may offer some benefits. The ionizing step can be driven with high intensity from a fixed-frequency pulsed laser, such as a frequency-doubled YAG laser. If the first step is driven with a narrow-linewidth cw laser at low intensity, rotational resolution should be achievable.
###### Acknowledgements.
We gratefully acknowledge support from NSF and the University of Connecticut Research Foundation. We thank Chad Orzel for making the pulsed dye laser available, and Ye Huang and Hyewon Kim for laboratory assistance.
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# Model Theory and the 𝐴𝑑𝑆/𝐶𝐹𝑇 correspondence.
## I Introduction and Motivation
The $`AdS/CFT`$ correspondence has been the subject to many substantial checks since its discovery by Maldacena maldacena1997 . The best known example of the correspondence, namely superstring on $`AdS_5\times S^5`$ background versus superconformal Yang Mills $`SU(N)`$ on 4-dimensional Minkowski space (with 4 supersymmetries), is especially well known and explicitly described witten1998 .
The correspondence expresses the fact that two completely different theories, the theory of (quantum) gravity - sustring in 10 dimensions and superconformal pure YM theory on 4-dimensional flat space, hence a theory without gravity, are completly equivalent to physical predictions they provide. Moreover, the correspondence correlates the weak - strong (perturbative - nonperturbative) coupling limits of the theories.
On the YM side of the correspondence a lot of effort went into obtaining less supersymmetric and nonconformal YM theory which would be as much realistic as possible. The goal would be the 4-dimensional QCD-like theory (or the Standard Model of particle physics) as completely dual to the sustring theory Cvetic2005 . This fascinating aim, however, has not been achieved yet. There exist many proposals how to proceed toward it. The purpose of this paper is to present arguments that some new mathematical tools can be relevant for such purposes. The tools in question are exotic smooth differential structures on the topologically trivial $`^4`$. However, one should refer to the formal mathematical objects in perspective established by the model-theoretic paradigm rather than ascribe to the absolute classical approach where various mathematical tools are placed in the absolute ”Newton-like classical” space, and governed by the ever present absolute classical logic.
We attempt to explain briefly reasons that model theory become relevant for the Maldacena duality or quantum gravity in general. First, the Maldacena duality, stating the equivalence of two different theories in different dimensions, is perfectly suited for applying the tools of category theory. Categories can carry their internal logic. Possibly, models of formal theories, when taken in general categories, can guide us how to describe the duality.
Next, the attempts to build a theory of QG are faced with extremly high energies, several orders higher than in Standard Model (SM) of elementary particles or in classical general relativity (GR). The proposed theories describing these limits, the theories of QG, rely heavily on the mathematical apparatus rather than experimental data. In such a case the theories can be especially sensitive to the degree of formalization of mathematics involved. If we had experimental data concerning the QG regime, it would be enough to propose the mathematical models explaining the results. In the opposite situation, which we face with, the mathematics with its delicate and sometimes subtle problems can be decissive for the successful formulation of QG. All the more that QG unified with other interactions is sometimes referred to as the ”theory of everything”, in the sense that it should deal with any physically valid situation. What we propose is to take into account also the model theoretic aspects of formal mathematical tools involved in the formulation of QG. This is connected with the necessarily increasing degree of formalization of the theory, and hence more and more formal approach requires referring to the nonuniqueness of formal models of the theories and languages used. Thus, not only objects derived from mathematics but also formal languages and their models can become valid objects for the full theory of QG (see also benioff1999 ; benioff2002 ; benioff2003 ).
Finally, the mathematical tools involved in QG should be applicable in all the scales of energy, from big bang, black holes to the everyday’s scales of energy. The mathematical formalism seems to be extremely invariant fixing the absolute point of view of the ”observer” with his absolute mathematical tools, like classical logic and real or natural numbers, that have to maintain relevance in all those extremly different physical regimes. It would be quite a surprise if there were no need to modify some tools when trying to describe such extremly different limits properly. All the more that mathematics itself can menage the objects properly if modified at very fundamental level, provided that the theories are sufficiently formal. The suitable tools derive from model theory.
Thus, our aims are twofold. First, we try to collect arguments that indeed the absoluteness of classical logic organizing the theory at the metalevel must be substantially weakened. The logic will then shift to intuitionism. Similarly, natural and real numbers can be considered as varying objects important for the description of QG. Both may be properly grasped by referring to some tools from model theory. Second, in the limit of varying fundamental objects as above, one can still refer to a smooth structure supporting the changes. This is exotic smooth differential structure on topological $`^4`$ and such a structure can be relevant for the generelized Maldacena duality.
Let us begin with some examination of the ”axiom” stating that the physical theory, if formulated mathematically, should be necessarily organized by classical logic at its metalevel, i.e. at the user level of this theory.
First, the coordinate spacetime frames that the theory referes to at some of its limits are usually understood physically as determined by some measuring devices like physical rods and watches. These, in turn, when described mathematically are usually understood as segments of real numbers. It is straightforward that the real numbers associated with the rods and watches have all the properties of the formal real numbers described by some sets of formal axioms. It is also known that the real axis is uniquely described in the second order predicate language. Thus, the rods considered formally require the reference to the second order properties. This means that some rich part of Zermelo-Fraenkel set theory with the axiom of choice (ZFC) is assumed to be valid. The question is: does the use of our physical measuring objects really justify it, and if not, what alternative do we have, if such alternative exists in mathematics at all?
Let us assume at the beginning that we have rods as physical objects but we refer directly to the points on the idealized 1-dimensional rod and we do not want to refer to the subsets of the points of this rod in the formal sense of ZFC. This seems to be physically natural assumption. Thus, our formal language, allowing for the above presentation of the rod, should be the first order predicate language with the names-symbols for every point of the one-dimensional ideal rod. In this way we have the correspondence between points of the rods and real numbers in some segment of reals. But we do not extend it over to the corespondence between subsets in the sense that we avoid stating the properties which require the higher order language.
Now, according to classical model theory, any theory in the first order predicate language has infinitely many different models of any cardinality greater than or equal to the cardinality of the language. Thus in order to retrieve uniqueness we can follow two possibilities: one is to raise the order of formal languages. But this is exactly what we have attempted to avoid. The second is to built a theory which would be invariant in a proper sense with respect to some class of the first order models. The latter is more favorable since we can maintain more physical and direct picture of measurements using rods – watches technology, and we can follow the Einsteinian idea of general relativity as if extended over nonstandard models of flat open regions of 4-dimensional spacetime.
Besides, the analogy between local frames and Boolean-valued models of ZFC in the context of QM was explicitly stated already in the early papers by Davis davis and Takeuti takeuti1978 . Then, the ideas where put forward in the context of arbitrary toposes with natural numbers object by Bell bell1986 .
Now, let us take the first order models of the theory of real numbers, i.e. all the first order properties of reals should hold in the models. Let us generate the models by the ultrafilter construction modulo some nonprincipal ultrafilter. It is known that these models are elementarilly equivalent to the standard models of reals KeislerChang . This means that all first order properties of real numbers are the same for both, standard and nonstandard models. One needs the higher order properties of real numbers to distinguish the models. However, without referring to the higher order languages, one can still attempt to formulate a physical theory which would be invariant with respect to the choice of the models of reals.
To distinguish between the models one has to refer to the classical and formal metalevel, which is ever present in an unchanged and absolute way. This is the point which requires relativization similar to the one performed in special or general relativity regarding the absolute ether or space and time.
However, this relativization is essentially involved in the foundations of mathematics and, as we will see, cannot be considered as merely physically valid concept.
Let us attempt to recover the information about the above relativization which is present already in the classical approach to manifolds and their smooth structures, and the relevance of both to background dependent or independent physical theories. The theories in question are sustring theory, or GR and loop QG respectively.
In a background dependent theory one referes to the structure of the underlying manifold (metric, local patches); in the case of sustring this is the structure of 10 (11) dimensional manifold with the fixed Minkowski metric. There are well-defined points of the manifold which one can refer to. The manifold’s structure derives, in fact, from the set-like structure. Besides, the metric allows one to measure distances between the points.
In a background independent theory we have metric as a dynamical variable and the theory should be diffeomorphism invariant with respect to the 4-dimensional diffeomorphisms. The point-like structure is not important any longer in the sense that the theory should produce quantities completely independent on the points of the manifold or distances between these. However, the quantities are sensitive to the smooth structure of the underlying manifold. Thus the description of the smooth invariants, hence the smooth structure, should not refer to the point-like distances.
The natural question can be asked: are there smooth structures defined such that the reference to the point-like level of the underlying manifold is limited in principle, with the exception of general knowledge of the topological type of the manifold?
Let us observe that: if we allow for the changes of the models of SET (ZFC), as defining the reals building the manifold, the changes can cause the limiting access to the point-like structure of the manifold due to the smearing of the reals between the models.
However, in principle, we should always recover the unique point-like description by the choice of sufficiently high classical metalevel where the models can be recognized and distinguished. Thus the change of the models of ZFC affects merely the formal presentation of the set-like point structure agreeing with its topological and smooth structures. They vary altogether along with the changes of the models of ZFC.
However, what if the change in the models of ZFC occurs but a smooth structure is fixed and, moreover, the changes of the models are essential for the description of the smooth structure?
In that case, indeed, the smooth structure cannot refer uniquely to the point-like SET structure, and this unrespecting of the level of points is essential for the description of the smooth structure. However, it would be enough to raise the metalevel to the ambient classical one such that one can distinguish the models and hence the shift in reals defining the manifold can be detected as model-dependent. Thus, the corresponding smooth structure would not be defined on a manifold with $`^n`$ local patches but rather the structure varies again from one model to the other.
However, raising the metalevel cannot be controlled formally in a full degree. Thus, it is not obvious that the raising the metalevel and attempting to keep it sufficiently formal give the classical result. Besides, in the complex theory – metatheory there are always present some informal ingredients giving the limitations for formalization Vaananen (see Appendices A and B).
Thus, attempting the further formalization of the complex theory – metatheory, the intuitionistic logic appears at the metalevel rather than the classical one krol2004 . Moreover, the lack of formal tools makes that the different formal models of basic mathematical structures, like reals or naturals, can happen to be indistinguishable. The models exist formally, although limitations at the metalevel make them indistinguishable. The attempt to formalize fully the metalevel results in the intuitionistic logic behind mathematical structures. The weakening of the logic gives the indistinguishability of the models.
The emergence of intuitionistic logic is evident when we start to approach arbitrary smooth real manifold literally by the symbols of some formal language. Then open subsets become important. The subsets generate topology as well. To extend this kind of description over all smooth manifolds from the category of these and smooth maps between them, one has to deal with the intuitionistic logic of the smooth topos. The topos is Basel topos which was constructed in ReyesMoerdijk .
The symbols of the formal language respecting the open covers of the manifolds appear to be interpreted already in Basel topos. This can be understood as generic situation for the emergence of the intuitionistic logic at the metalevel when formalized. As the consequence two different models of reals appear which cannot be formally distinguished.
We refer the reader to the Appendices A and B for more detailed explanation of these matters.
In such a way we can be ensured about the existence of some smooth differentiable structure on some smooth manifold which does not derive from the fixed SET structure of the manifold. This structure is identified with exotic smooth $`R^4`$ krol2004 . The point-like level of this manifold dynamically changes between some models of ZFC, and because of that the points level does not generate the structure. This is rather the dynamics of the changes of the models of ZFC. Surprisingly, such a structure is better suited for the expressing the background independent quantities, as required by GR or QG, then the referring to the standard smooth structures, which are tightly connected with the point-like structures of the smooth manifolds. Besides, the invariants of the exotic $`R^4`$’s might be generated by respecting the changing models of ZFC, and these invariants can be useful in the background independent formulation of QG krol2004 .
In the next three paragraphs we present the reasons for considering exotic smooth $`R^4`$’s in the context of GR, QG and the Maldacena duality. The model theoretic constructions, essential for establishing the connection, are discussed in Appendices A and B.
## II Why exotic smooth structures in QG?
Recently Pfeiffer pfeiffer2004 has presented some arguments that exotic smooth structures, on compact as well as on noncompact 4-manifolds, should be enormously important for quantum gravity (QG) formulated in a path integral form, i.e. as a sum over geometries on compact, smooth 4-manifolds with boundaries. The relevance of exotic smooth structures on noncompact $`^4`$ for QG formulated as a background independent theory was conjectured already in krol2002 ; krol2004 . In this and subsequent paragraphs we attempt to go further in explaining the relevance of exotic smoothness in dimension 4 for QG program. In fact, we argue that exotic smooth $`R^4`$’s can play a significant role for QG, which is similar to the role played by standard $`^4`$ in classical general relativity because of its local behaviour. Coordinate axes are considered usually as the set of real numbers. Nevertheless, such entities can be formally perceived as the models of the theory of real fields. The models are unique only in the higher orders of the formal language (in fact, the second order is sufficient). This interplay between models and formal languages gives some kind of an additional degree of freedom to a theory.
The success in generating nontrivial combinatorial 4 manifold invariants by QG, could prove essential for the pure differential geometric task of classifying and distinguishing different nondiffeomorphic smooth structures on 4-manifolds via their PL-structures. This was partially achieved by Mackaay state sums or Crane-Yetter smooth invariants mackaay1 ; mackaay2 ; crane1997 . Donaldson and Seiberg-Witten smooth invariants are known to distinguish exotic smooth structures on compact closed 4-manifolds, but it is also known that these invariants are not refined enough to distinguish between all different nondiffeomorphic smooth structures on some 4-manifolds (see stern1999 for the case of non simply connected 4-manifolds). Thus, to find 4-invariants which would be powerful enough and combinatorial is a challenge for differential geometry. Moreover, in dimension four one can have infinite continuum many nondiffeomorphic different smooth structures on noncompact 4-manifolds. They also have their PL structures associated uniquely, so to calculate suitable invariants in this case is again challenging taylor1997 .
Since the ”early” days of the development of topological quantum field theories (TQFT) atiyah1988 ; atiyah1989 ; witten1988 it has been realized that the connection of physical quantum theories with some subtle mathematical questions regarding smoothness of 4-manifolds is very deep and nontrivial. Witten witten1988 ; witten1994 was able to show explicitly how one can obtain Donaldson polynomial invariants as correlation functions in certain twisted supersymmetric, $`𝒩=`$ 1 Yang-Mills theory on compact smooth 4-manifolds. This Yang-Mills theory possesses some features of the realistic QCD. Namely it is formulated in 4 dimensions, with an asymptotic freedom and chiral symmetry breaking, and with a dynamically generated mass gap. Donaldson polynomial invariants are smooth 4 compact, oriented manifold invariants invented by the nontrivial use of the Yang-Mills theory technics donaldson1990 . They should distinguish between different exotic smooth structures of compact closed oriented (simply connected) 4-manifolds. Witten and Seiberg proposed to calculate some other smooth 4-invariants, the Seiberg-Witten invariants, by the inherent use of physical ideas dealing with some dualities of the theories involved witten1994a . The relation between Seiberg-Witten and the Donaldson invariants is mathematically nontrivial and is a subject of some profound conjectures donaldson1990 .
The involvement of ideas from physics in the creation of the invariants is not accidental. It is higly essential and the relations are crucial for both fields: physics and mathematics. This is a vast field of investigations by itself and in this paper it is only touched upon. The main lesson is that there is crucial involvement of 4-smoothness in some physical field theories and appealing to the physical concepts is indispensable when pure mathematical questions, regarding 4-smoothness, appear. This should not be given up only on account of the presence of some technical complications.
In $`3+1`$ quantum gravity the connection is not less profound. GR needs 4-dimensional differential structure of 4-manifolds in the very intrinsic way. Moreover, formal partition functions as calculated in path integral formulation of QG (if it is formulated) have to be smooth invariants of the underlying compact, closed 4-manifolds pfeiffer2004 . Moreover, formal partition function of quantum general relativity as computed by purely combinatorial means is an invariant of PL 4-manifolds. PL-invariance is guranteed by an invariance with respect to the finite sequences of the Pachner moves pfeiffer2004 . For any smooth manifold of dimension $`d6`$ any smooth structure determines uniquely a combinatorial PL-structure and conversly, any PL-structure determines uniquely a smooth structure of d-manifold gompfStipshitz1999 . Thus, the partition function is automatically an invariant of the corresponding 4-dimensional smooth manifold. These PL-invariant partition functions generally take the form of state-sums and can be schematically represented as follows:
$$Z=\underset{\{colourings\}}{}\underset{\{simplices\}}{}(amplitudes).$$
(1)
Any colouring labels the faces of simplices from the triangulation of the manifold. Amplitudes are numbers-traces associated with any such labeling pfeiffer2004 . All mastery with construction of suitable invariants is the appropriate choice of the category where the sets of colourings derive from, and such that the invariance with respect to the Pachner moves holds. If so, one has PL, hence diffeomorphism invariance of the (1). We do not review here the existing invariants in dimensions 3 or 4, since there is a vast literature on this subject pfeiffer2004 , however the general observation can be made, namely while shifting from the lower to the higher dimension suitable higher n-categories have to be involved. This became clear especially after the publications of Baez and others in the context of TQFT, where a canonical functor was considered, the one from the category of $`n`$-cobordisms to the $`n`$-category of $`n`$-vector spaces baez1995 ; baez2004 . Crane and Yetter proposed to build 4-manifold invariants based on the category of finite dimensional representations of the quantum group $`U_q(sl(2))`$ where $`q`$ is a principal $`4r`$-th root of unity. Then, it became evident that it is possible to construct 4-manifold invariants out of some special 2-categories. Indeed, a detailed construction was performed by Mackaay mackaay1 . The 2-category he referred to was some kind of spherical 2-category (in the simplest variant it is a semi-strict 2-category of 2-Hilbert spaces which was the completely coordinatized version of 2Hilb category as introduced by Baez baez1995 ). The highly categorical context, when 4-dimensional smoothness is considered, is also apparent in the model-theoretic approach to the subject. In the last section We will make some comments on the possible way of generating 4-smooth invariants via model-theoretic means.
Following Pfeiffer pfeiffer2004 , from the point of view of physics and QG, the existence of state-sum invariant of PL 4-manifolds would be of special interest. The invariant in the degenerate limit should yield the path integral quantization of $`d=3+1`$ pure general relativity. The degenerate limit is understood by analogy with the limit of $`2+1`$ Turaev-Viro state sum invariant which yields the degenerate Ponzano-Regge model. The latter deals with the background independent quantization of pure GR in 3-dimensions with the vanishing cosmological constant and with Riemanian signature of the 3-metric. Degeneracy of the limit means that the limiting partition function yielded from the Turaev-Viro invariant is divergent, however, after removing the bad terms, it is still invariant with respect to the Pachner moves pfeiffer2004 . We will understand that in the case of $`3+1`$ invariants the ”regularization” of the degenerate limit can be connected with special model-theoretic properties of exotic $`R^4`$’s.
Moreover, in the context of sustring theory it was shown strominger1993 that certain scattering amplitudes of fivebranes solitons and axions in the heterotic string theory are proportional to the Donaldson polynomials. Thus, 4-dimensional exotic smoothness is rooted in the formalism of string theory as well.
## III Deformed Special Relativity program and exotic smooth $`R^4`$’s
The possibility that some exotic $`R^4`$’s can generate QG effects krol2004 leads to the reverse possibility that some QG effects could be compensated by a suitable choice of some exotic $`R^4`$’s. In this section we attempt to analyze the range of the validity of this later question by, first, dealing with the special limit of QG, the one which is explored in doubly special relativity program, and then, considering more general situations of full QG.
Deformed or Doubly Special Relativity (DSR) program was invented amelino2001 ; amelino2002 <sup>2</sup><sup>2</sup>2I thank professor Giovanni Amelino-Camelia for letting me know about his papers as a very direct and attractive attempt to find an extension of Lorentz or Poincaré algebra such that it would respect also a fundamental role played by Planck scale of length in the quantum gravity regime glikman2003 ; oriti2004 . It was assumed that the length should be invariant with respect to the extended algebra as another, next to the speed of light, invariant quantity. Then, it became clear that this requirement led to the necessity of considering noncommuting coordinates in the tangent and base spaces. In fact, spacetime coordinates are generated by translation generators in the energy-momentum space glikman2003 . The generators do not commute and they can be realized as the generators of the remaining part while the decomposition of the symmetries $`SO(4,1)`$ of 4-dimensional de Sitter space into Lorentz group $`SO(3,1)`$ and the remainder is performed.
It is known that $`SO(4,1)`$ being symmetry group of the 4-dimensional de Sitter space is also the conformal group of the Euclidean 3-space $`^3`$ petersen1999 . $`SO(4,1)`$ is a subgroup of $`SO(5,1)`$ \- the isometry group of 5-dimensional de Sitter space. $`SO(5,1)`$, in turn, is the conformal group of the 4-dimensional Euclidean space $`^4`$; hence, the 4 dilations are included in the later group. Besides, 4-dimensional Lorentz group is present in the decomposition of $`SO(4,1)`$. In other words Lorentz group is the subgroup of the 4-dimensional conformal group. In fact, n-dim Poincaré group has $`\frac{1}{2}n(n+1)`$ generators while the conformal group in n-dimensions has $`\frac{1}{2}(n+1)(n+2)`$ generators. The difference is $`n+1`$ generators which $`n`$ of them generate special conformal transformations and remaining one parameter corresponds to the 4-dilations.
It was conjectured in krol2004b that some standard (with respect to the standard smooth structure on $`^4`$) contractions and, possibly, rotations of the 4-ball where some exotic $`R^4`$ is placed give, as the result of such localization, some non commuting ”momentum”coordinates corresponding to some noncommutative space. In fact, any arbitrarily small contraction makes the exotic smooth structure being deformed such that it cannot be embedded in the ambient standard structure extending the exotic one gompfStipshitz1999 .
We can synthesize the whole idea by assuming that some 4-dilations of some exotic 4-space give the noncommutative space, which in some special limit gives the noncommutative 4-space exactly as in the case of DSR noncommutativity i.e. generated by four-dimensional de Sitter momentum space. The limit of QG where DSR approach is valid is the ”flat space semi-classical limit of $`3+1`$ quantum gravity”. This means that one takes Newton constant $`G0`$ and Planck constant $`h0`$ but such that their common ratio $`\sqrt{\frac{G}{h}}`$ remains finite. This finite limiting ratio is in fact Planck length and the vanishing $`G`$ limit of GR is flat spacetime solution of topological field theory; taking $`h0`$ is the weak coupling limit of QG. The finite fraction $`lim_{G,h0}\sqrt{\frac{G}{h}}=\kappa ^1`$ corresponds to $`\kappa `$ which is exactly the radius of the 4-de Sitter space defined in the 5-dimensional Euclidean space as a surface: $`\kappa ^2=\eta _0^2+\eta _1^2+\eta _2^2+\eta _3^2+\eta _4^2`$ where $`\eta _0,\eta _1,\eta _2,\eta _3,\eta _4`$ are coordinates in the Euklidean space glikman2003 .
Although the relation between some contracted exotic $`R^4`$ and noncommutative spaces was observed to exist, the relation does not decide the exact shape of the noncommutative space. However, it is plausible that the class of noncommutative spaces associated with contracted exotic $`R^4`$’s appears to be wide enough to state the following general conjecture:
A1. In the flat space semi-classical limit of $`3+1`$ quantum gravity glikman2003 one reaches the noncommutative DSR space, as connected with some reference frame of the momentum 4-de Sitter space, which is exactly, possibly deformed, the noncommutative space generated by some contracted exotic $`R^4`$ krol2004b .
or
A2. Any quantum gravitational effect in the flat space semi-classical limit of QG can be locally compensated by a suitable choice of some exotic smooth $`R^4`$.
Let us develop some arguments that would be relevant to another stronger hypothesis, namely the one stating that exotic $`R^4`$’s can be involved more essentially in a formulation of full QG than only in the limit as above.
* First, exotic $`R^4`$’s are tools suitable for the detection or measurement of exotic 4-dimensional smoothness in the case of compact 4-manifolds. This is precisely the way followed by Taylor taylor1997 in generating the smooth invariants of compact or noncompact 4-manifolds. However, compact exotic smoothness is essential for QG, at least in the path integral formulation pfeiffer2004 , hence, also some exotic noncompact $`R^4`$’s should be considered important there.
* Second, spin networks or state sums are natural states of loop QG which is a background independent formulation of QG in 4 dimensions baez1995 . Besides, natural candidates for smooth PL invariants of compact 4-manifolds are state sum invariants mackaay1 . In the knowledge that exotic smooth $`R^4`$’s are related to the compact smoothness as in the previous point, one comprehends that exotic $`R^4`$’s should be important for a background independent formulation of QG.
* Third, considering exotic world volumes of $`D3`$-brabnes in sustring theory, exotic $`R^4`$’s may have a meaning in the perturbative sustring description of QG krol2004 . This can be further exploited in the context of the Maldacena conjecture krol2004 .
* Fourth, exotic $`R^4`$’s smooth invariants, if calculated by some model theoretic means (MTSD networks), can compansate infinities or divergencies (see section V). This, in turn, can help to understand perturbative, though, non renormalized, quantum gravity.
Moreover, let us examine some general arguments tackling symmetries involved in the DSR and contracted exotic $`R^4`$’s. The standard contractions of exotic $`R^4`$’s require reference to the 4 dimensional contractions which are elements of the conformal group of 4-space $`^4`$. This last group is $`SO(5,1)`$ or in the case of Minkowski’s metric the conformal group is $`SO(4,2)`$. In turn, $`SO(4,2)`$ is the isometry group of anti-de Sitter 5-space, namely $`AdS_5`$. Thus, to compensate properly QG effects in the flat space semi-classical limit of QG in terms of exotic self-dual smooth $`R^4`$’s one should deal with a broader group than $`SO(4,1)`$. This broader group has to contain 4-dilations of Euclidean or Minkowski space; this might be $`SO(5,1)`$ or $`SO(4,2)`$. Any field theory, respecting these groups of symmetries and being formulated on flat, say Minkowski $`^4`$, while the standard smoothness of the underlaying $`^4`$ patch is changed to the exotic one, should contain a reference to the quantum gravitational regime as well. If the field theory on standard Minkowski 4-space respects the $`SO(4,2)`$ symmetry, this theory is clearly conformally invariant. Now if one switches the smooth structure to the exotic one, the field theory will not be conformally invariant any longer, since in the regime of suitable contractions the noncommutative (quantum gravitational) effects appear, or in terms of exotic smooth structure the exotic metric is not in the conformal class of the Minkowski metric. Moreover, the exotic metric need not be standard smooth one as well as the Minkowski metric is not a smooth metric in the exotic structure.
Now if one considers the basic case of the Maldacena conjecture, namely the case of sustrings on $`AdS_5\times S^5`$ dual to the $`SU(N)`$ Yang-Mills superconformal with $`𝒩=`$ 4 and $`N`$ large, he is faced with the possibility that the conformal invariance of the theory should be broken due to the possible intervention of exotic $`R^4`$ structures krol2004 . Let us observe that the $`AdS_5`$ is just the space whose isometry group supports conformal transformations of 4-Minkowski space and this was the initial test of the duality. Moreover, the breaking of conformal invariance in this context should be done in order to have the semi-classical limit of 4 dim QG well adapted.
Thus, we can see that the involvement of exotic smooth $`R^4`$’s is not only limited to the flat space semi-classical limit of $`3+1`$ QG, but it may play more profound role in full QG. The observations encourage us to formulate the following conjecture:
A3. Any quantum gravitational effect in 4-dimensions, perturbative or non-perturbative, can be compensated locally by a suitable choice of some exotic $`R^4`$.
The family of exotic $`R^4`$’s emerges, such that $`R^4`$’s compensate QG effects locally. This family remains in some analogy to differentiable manifolds and its coordinate local $`^4`$ patches, as in the case of classical gravity. However, this family does not constitute any cover of differentiable manifold in the strict, classical sense.
We have presented only some heuristic, general arguments for the above conjecture, however, one can become more convinced about its validity by a more careful examination of model-theoretic self-dual exotic smooth structures on $`^4`$. We will see that the model-theoretic approach to exotic smoothness resembles to a certain degree the philosophy behind the general relativity approach to the classical gravity. Certainly, the very meaning of this conjecture requires explanation. At our disposal these are two possible scenarios of studying the exotic $`R^4`$’s and their meaning in physics. One is external, based on general relations to some other well-known concepts or tools and the second internal, and it attempts to grasp exoticness from the point of view of its intrinsic nonclassical perspective (see Apendices A and B). It is strange and unexpected that the latter is more suitable for exhibiting the physical meaning of the exotic structures. On the other hand an impasse in a direct, analytical approach to exotic smoothness in dimension four caused the latter to emerge krol2002 ; krol2004 ; krol2004b .
## IV The general features of exotic $`R^4`$’s indicating that they are well suited for General Relativity and Quantum Gravity.
Before one attempts to obtain concrete calculations concerning exotic $`R^4`$’s and their connection with QG or GR it is worth examining at the problem from a more general perspective and to collect some general indications showing that the connection should be studied.
* First, the conjecture III.A3 resembles much the Einsteinian equivalence principle (EP) in GR. Let us have a closer look at this.
At least two versions of classical EP exist in general relativity weinberg . EP says that at every spacetime point in an arbitrary gravitational field it is always possible to choose a ”locally inertial coordinate system” such that, within a sufficiently small region of the point, the laws of nature take the same form as in unaccelerated Cartesian coordinate system in the absense of gravitation.
The weak EP holds when in the above statement ”the laws” means ”the laws of motion of freely falling particles”. Strong EP holds if we understand ”the laws of nature” as ”all laws of nature”.
EP is connected directly to general covariance or the general feature of Lorentzian manifolds (of any manifold in fact) where it is always possible to take a locally trivial coordinate patch. A general diffeomorphisms generate gravity via the metric and Christoffel connection which in turn results in the concept of covariant derivative. The information deriving from the covariant derivative characterizes the tangent bundle of the manifold in global.
Now let us reformulate EP as follows: In spacetime it is always possible to choose locally standard $`^4`$ patch in which no gravitational effects are observed, or
A3’. The standard $`^4`$ patch can be chosen such that this compensates locally all classical effects of gravity.
Now the resemblance to the quantum gravitational counterpart III.A3 is striking. Both statements are deeply rooted in the differential geometrical properties of four-manifolds.
* Second, the general invariance, or coordinate independence, with respect to 4-diffeomorphisms required by GR assures us that every decomposition of spacetime into 3 space directions and time direction is a violation of the invariance. But when this happens it can be recovered just by summing over all decompositions or by showing that the result does not depend on the specific decomposition performed. However, the general invariance or the obstruction with decomposing $`^4`$ into $`\times \times \times `$ is such a fundemental feature of GR that while reconciliation with QM is considered this has to be carefully analysed. Let us note that the above decomposition agrees with the topological understanding of coordinate axes and their smooth product is understood as compatible with the topological product. Such an implicit assumption enforces one to deal with the standard smoothness, if any, of $`^4`$. Moreover, another implicit assumption is made. This is the order of decomposing into coordinates and recovering the smooth structure on $`^4`$. This order is assumed to be irrelevant and the above actions can be performed anytime and they always result in the standard smooth structure. The result does not depend on these factors and one has always control of the compatibility of the smooth and topological structures. As a result, GR, though generally invariant and coordinate independent, is built as if it were always equivalent to the coordinate dependent formulation or, as if the passage between the two were always in reach. This low level implicit assumption is however not made in the GR spirit.
Let us point out what exotic $`R^4`$’s can change regarding this point. Similarily as in GR, any choice of global smooth coordinates in the accordance with $`\times \times \times `$ destroys the exotic smoothness. However, its analytical restoration is not available at present. The ease with restoration of general invariance in the case of GR compatible theories as mentioned above, comparing with the dificulty, or even untractability with the explicit restoration of the exotic smoothness, should be an indicator that possibly something is overlooked in the GR/QM relations. Besides, there are arguments for the fact that one should take into account 4-dimensional exotic smoothness when QG is considered pfeiffer2004 .
On the other hand, the difficulty with the recovering of the exoticness from the standard topological coordinates shows that the topological axes are very deeply and canonically entangled in the exotic $`R^4`$. This is the feature required by GR. A theory on exotic $`R^4`$ respects the fundemental GR requirement. Moreover, contracted exotic $`R^4`$’s can interpret canonical noncommutative relations as in QM krol2004b and hence, heuristically, they appear to be perfect tools for unifying QM and GR.
What is more, it follows from krol2004b that model-theoretic self-dual exotic $`R^4`$’s exhibit noncommutativity when one settles fixed model-theoretic standardness and smooth standardness of coordinates. The last one is compatible with the topological $`^4`$ structure. That is why fixing that the topological and smooth structures agree, results in the fact that the QM-like behaviour is not any longer present. From such a perspective, it is even essential to take into account exotic $`R^4`$’s in order to reconcile QM and GR.
* Third, it was shown by Asselmeyer asselmeyer that exotic smooth structures on some compact 4-manifolds generate distribution-like sources of gravity. Besides, Brans brans1994 conjectured that exotic smooth $`R^4`$ can act as sources of gravitational field. It was Sładkowski sladkowski2001 showed that, indeed, some exotic $`R^4`$’s can generate highly nontrivial solutions of Einstein equations even in the empty but exotic $`R^4`$. The sources of gravity are exhibited while the exotic smooth structure is seen from the standard one, so this shift of the structures switches on the classical sources of gravity. In fact, it was proposed krol2002 ; krol2004 that such a shift can switch on some QG effects. This is also expressed in III.A3.
* Forth, when the background independence issue is considered in QG, the canonical role played by exotic $`R^4`$’s can be noticed krol2004 . It relies on the model-theoretic approach to exotic smoothness. The point is that the background manifold of the theory is frozen unless we take the model-theoretic dimension of the objects. Owing to this new perspective the manifold can vary. This dimension respects the model extensions of standard real and natural numbers. Such a situation generates model-theoretic self-dual exotic $`R^4`$’s and by Brans conjecture they can act as external sources of gravitational field. Consequently, model-theoretic background independence generates exotic smooth 4-structures and this has its gravitational counterpart.
* Fifth, assume we have QG formulated such that it respects somehow 4-dimensional exotic smoothness (e.g. path integral formulation). The classical limit in this case, at least locally, should correspond to the relation between exotic $`R^4`$ and standard smooth $`^4`$. The relation between the manifolds, at the topological level, is trivial, namely, both manifolds are topologically identical. This $`^4`$, while standard smooth, is just the main ingredient of classical GR in the local behaviour.
## V Exotic renormalization in AdS/CFT
Owing to the recent growth of interest in the geometry–field theory limit of the AdS/CFT correspondence, many explicit formulas have been worked out in favour of this holographic like relation skenderis2004 ; deharo2001 . We attempt to connect the so-called holographic renormalization procedures with the possible interventions of exotic smooth structures on $`^4`$. The analysis makes an essential use of the model-theoretic self-duality of exotic $`R^4`$’s.
The general strategy is as follows:
* Localize divergent quantities, as usual, with respect to some standard smooth and model-theoretic standard (MT standard) patch $`^4`$.
* Shift the patch to the one which is Robinsonnian nonstandard robinson1964 and the divergencies can be considered as nonstandard reals or nonstandard expressions in general.
* Change the model-theoretic environment to the intuitionistic one where the nonstandard reals or expressions can be considered as smooth-finite ($`s`$-finite) ReyesMoerdijk in the sense of the internal logic of Basel topos $``$ (see section (B)).
* Refer to the fact that there may exist model-theoretic self-dual (MTSD) exotic $`R^4`$ which can interpret all the above steps krol2004b .
* Attempt to consider a theory which is invariant with respect to the exotic diffeomorphisms of the exotic $`R^4`$ rather than to the standard ones krol2004b .
Such a theory has been built in a kind of renormalization procedure automatically, in the sense that the shift $`infinitefinite`$ is legitimate with respect to the general symmetry of the theory. Moreover, owing to the Brans conjecture brans1994 and the connection of exotic $`R^4`$’s and QG, the exotic MTSD $`R^4`$ suggests a connection with gravity (in both classical and quantum regimes) while the above renormalization scheme is been considered.
Since many explicit expressions heve been worked out in the holographic renormalization program we can also try to be more specific with our model-theoretic renormalization approach, at least in the geometry - field theory correspondence.
The conformal boundary of $`AdS_5`$ is $`S^4`$ in the Euclidean picture of the $`AdS_5`$ space witten1998 (This boundary can be considered as $`^4\{\mathrm{}\}`$). Potentially, $`^4`$ can carry an infinite number of different smooth structures but after the one–point compactification they might be unique again (up to the smooth 4-dim Poincaré conjecture which is still open gompfStipshitz1999 ).
Now, let us consider $`S^4`$ as the one-point compactification of $`^4`$, but in the sense of Robinson’s nonstandard analysis robinson1964 ; this compactification is the quotient of the nonstandard exstension, $`{}_{}{}^{}𝐑_{}^{\mathrm{𝟒}}`$, by some equivalence relation (nan1985 , p. 159).
For this reason divergent functions on $`S^4`$ might be considered as nonstandard functions on $`{}_{}{}^{}𝐑_{}^{\mathrm{𝟒}}`$ up to the equivalence relation defining the one point compactification.
Now, let us consider the gravity counterterms to the $`AdS_5`$ Einstein-Hilbert (EH) supergravity action, which have been calculated and presented in several papers (see e.g. skenderis2001 ; skenderis2002 ; skenderis2004 ). These counterterms are needed to compensate for the divergencies in the EH action. The action is formulated in the bulk which in this case is $`AdS_5`$. EH action supplemented by a surface term, on anti-de-Sitter $`n`$-dimensional space, can be written as
$$S_{\mathrm{bulk}}+S_{\mathrm{surf}}=\frac{1}{16\pi G}_{AdS}d^{n+1}x\sqrt{\stackrel{~}{g}}\left(R+\frac{n(n1)}{l^2}\right)\frac{1}{8\pi G}_{(AdS)}d^nx\sqrt{h}K$$
(2)
The bulk term is just EH action with the cosmological constant $`\mathrm{\Lambda }=\frac{n(n1)}{2l^2}`$, $`l`$ is the radius of the $`AdS`$ space, $`K`$ is the trace of the extrinsic curvature while embedding of the boundary $`(AdS)`$ in the $`AdS`$ is being considered. The boundary term guarantees that when field equations are calculated with respect to the above action (i.e. the variation with respect to the metric fixed at the boundary via its normal derivatives is calculated) they become Einstein equations emparan1999 . $`h`$ is the metric induced at the boundary from the bulk metric $`g`$ in the above sense. The action (2) is divergent, since the bulk is noncompact $`AdS`$ space and hence it has infinite volume, and the boundary term is divergent since the induced metric $`h`$ diverges on the boundary. Nevertheless, one can extract precise divergent counterterms. They are functionals of the boundary curvature and its derivatives only, and the counterterms depend on the radius $`l`$.
In the specific limit of the AdS/CFT correspondence we are interested in in the context of holographic renormalization, the classical supergravity theory in the asymptotic bulk $`AAdS_5\times S^5`$ is equivalent to the quantum field theory on the conformal 4-boundary. $`AAdS`$ is asymptotically anti-de Sitter 5-space skenderis2002 . Having the radial bulk coordinate $`\rho `$ such that $`\rho =0`$ corresponds to the asymptotic boundary, we can regularize the action taking surface corresponding to $`\rho =ϵ`$ skenderis2001a . Thus, the regularized action can be written as skenderis2002
$$S_{\mathrm{reg}}=\frac{1}{16\pi G}_{\rho 0}d^{4+1}x\sqrt{\stackrel{~}{g}}\left(R+\frac{6}{l^2}\right)\frac{1}{8\pi G}_{\rho =ϵ}d^4x\sqrt{h}K$$
(3)
Counterterms of the action in the case of $`AdS_n`$ space are
$`S_{\mathrm{ct}}={\displaystyle \frac{1}{8\pi G}}{\displaystyle _{(AdS)}}d^nx\sqrt{h}F(l,,)=`$ (4)
$`{\displaystyle \frac{1}{8\pi G}}{\displaystyle _{(AdS)}}d^nx\sqrt{h}\left[{\displaystyle \frac{3}{l}}+{\displaystyle \frac{l}{2(n2)}}+{\displaystyle \frac{l^3}{2(n4)(n2)^2}}\left(_{ab}^{ab}{\displaystyle \frac{n}{4(n1)}}^2\right)+\mathrm{}\right].`$
The counterterms corresponding to the regularized action on the $`\rho =ϵ`$ surface can be written as
$`S_{\mathrm{ct}}^ϵ[g_0]={\displaystyle \frac{1}{16\pi G}}{\displaystyle _{\rho =ϵ}}d^nx\sqrt{h}\left[2(1n)+{\displaystyle \frac{1}{n2}}{\displaystyle \frac{1}{(n4)(n2)^2}}\left(_{ab}^{ab}{\displaystyle \frac{n}{4(n1)}}^2\right)\mathrm{log}ϵa_{(n)}\right]`$
where $`\mathrm{log}ϵa_{(n)}`$ is the logarithmic term while the $`S_{\mathrm{reg}}`$ action is expanded in powers of $`ϵ`$ as skenderis2001a
$`S_{\mathrm{reg}}^ϵ={\displaystyle \frac{1}{16\pi G}}{\displaystyle d^nx\sqrt{g_{(0)}}\left(ϵ^{n/2}a_{(0)}+ϵ^{n/2+1}a_{(2)}+\mathrm{}+ϵ^1a_{(n2)}\mathrm{log}ϵa_{(n)}\right)}+𝒪(ϵ_{(0)})`$
and the metric in the neighbourhood of the boundary in the bulk is skenderis2002
$$ds^2=\stackrel{~}{g}_{\mu \nu }dx^\mu dx^\nu =\frac{1}{r^2}(dr^2+g_{ij}(x,r)dx^idx^j.$$
(5)
Now the boundary is located at $`r=0`$ and in the limit of $`r0`$ the metric has a smooth limit $`g_{(0)ij}`$. Hence, it can be written
$$g_{ij}(x,r)=g_{(0)ij}+rg_{(1)ij}+r^2g_{(2)ij}+r^3g_{(3)ij}+\mathrm{}.$$
(6)
This limiting, boundary metric $`g_0`$ was used in (V) and $`r`$ is the defining function of the boundary metric, namely it holds skenderis2002
$$g_{(0)}=r^2\stackrel{~}{g}_{|AdS}$$
(7)
and the bulk metric $`\stackrel{~}{g}`$ is evaluated at the boundary of $`AdS`$.
Introducing new coordinates $`\rho =r^2(=ϵ)`$ one can write skenderis2002
$`ds^2=\stackrel{~}{g}_{\mu \nu }dx^\mu dx^\nu ={\displaystyle \frac{d\rho ^2}{4\rho ^2}}+{\displaystyle \frac{1}{\rho }}g_{ij}(x,\rho )dx^idx^j`$ (8)
$`g(x,\rho )_{ij}=g_{(0)ij}+\rho ^{1/2}g_{(1)ij}+\rho g_{(2)ij}+\rho ^{3/2}g_{(3)ij}+\mathrm{}+\rho ^{n/2}g_{(n)ij}+\stackrel{~}{h}_{(n)ij}\rho ^{n/2}\mathrm{log}\rho +\mathrm{}.`$ (9)
Now subtracting all divergent terms $`S_{\mathrm{ct}}^ϵ[g_0]`$ from $`S_{\mathrm{reg}}^ϵ`$ in (V), while still staying on the regularizing surface $`\rho =ϵ`$ and performing $`ϵ0`$ limit in which renormalized action $`S_{\mathrm{ren}}[g_{(0)}]`$ takes form skenderis2002
$`S_{\mathrm{ren}}[g_{(0)}]=`$ (10)
$`\underset{ϵ0}{lim}\left(S_{\mathrm{reg}}{\displaystyle \frac{1}{16\pi G}}{\displaystyle _{\rho =ϵ}}\sqrt{h}\left[A\right]\right)`$
$`\mathrm{where}A=2(1n)+{\displaystyle \frac{1}{n2}}{\displaystyle \frac{1}{(n4)(n2)^2}}(_{ij}^{ij}{\displaystyle \frac{n}{4(n1)}}^2)\mathrm{log}ϵa_{(n)}.`$
Thus, all finite number of divergent terms in the action were expressed in terms of the boundary induced metric and then, when the divergent terms were subtracted, the renormalized action is finite and well-defined. Now one can produce various boundary $`n`$-point functions via field theory calculations skenderis2002 and they should precisely characterize 5-dimensional bulk geometry of the asymptotic anti-de Sitter space.
Now, following solodukhin1999 , let us approach the 5-dimensional noncompact bulk by the increasing sequence of compact embedded submanifolds with boundaries. Therefore, we have submanifolds $`M_\rho `$ with boundaries $`M_\rho `$ parametrized by $`\rho `$ such that $`M_\rho M`$ for large $`\rho `$. The action functional (2) as written on the $`M`$ should be approached by the sequence $`\{S_\rho \}`$ of the corresponding functionals. This means that divergencies of the limiting functional are also approached.
Now instead of taking the limit $`M_\rho M`$ and $`M_\rho M`$ and renormalizing the action via counterterms which are the functionals of the curvature invariants at the boundary as discussed above, let us take the limit modulo some nonprincipal ultrafilter $`U`$ on the infinite set of indices $`\{\rho \}\omega `$ nan1985 . As the result, we get two ingredients
* a nonstandard boundary $`{}_{}{}^{}M(\{M_{\rho _n}\}^n\mathrm{})/U`$, which locally can be identified with the nonstandard $`{}_{}{}^{}𝐑_{}^{\mathrm{𝟒}}`$,
* a nonstandard action containing some symbols which are nonstandard or which refer to some nonstandard real numbers.
The first ingredient can be considered as the generating nonstandard $`{}_{}{}^{}𝐑_{}^{\mathrm{𝟒}}`$ which was used in the nonstandard construction of the one-point compactification of $`^4`$ at the beginning of this section. The second ingredient can be worked out when one takes in the renormalizing limit $`ϵ0`$ in (10) or (V) just sequences of $`ϵ`$’s as corresponding to $`\rho `$’s. Then, on such a family of the expressions let us perform the modulo nonprincipal ultrafilter operation. Thus, the divergencies of the counterterms $`S_{\mathrm{ct}}^ϵ[g_0]`$ as in (V) correspond to nonstandard reals, when the model-theoretic limit (i.e. modulo the nonprincipal ultrafilter) is taken and when the nonstandard boundary is generated as well. Moreover, the limiting action now can be seen as the nonstandard functional defined on the nonstandard functions which, in turn, are defined locally on $`{}_{}{}^{}𝐑_{}^{\mathrm{𝟒}}`$ with the values in $`{}_{}{}^{}𝐑`$ or $`{}_{}{}^{}`$.
The main point of this model-theoretic construction is that having the sequence of the submanifolds and their boundaries which carry canonical order with respect to the increasing parameter $`\rho `$, one does not respect this order any longer and makes the identification of the boundaries modulo some ultrafilter. This new limiting procedure is just the one which replaces the old, order-like, one. However, there is still a connection with the former limiting procedure. This is clearly the MTSD property of exotic $`R^4`$’s. Having the appropriate MTSD exotic $`R^4`$, one can choose the $`^4`$ patch where ”standard” infinities of the action are detected or switched to the nonstandard $`{}_{}{}^{}𝐑_{}^{\mathrm{𝟒}}`$ patch where the infinities are still real, though nonstandard, numbers. The shift between the patches can be compensated in principle by some intuitionistic environment modelled by the inside of the Basel topos. The smooth structure which can afford all the changes is just model-theoretic self-dual exotic $`R^4`$. Hence, instead simply dropping the infinities, one can create a theory which can model canonically the renormalization procedure. The symmetries should respect MTSD exotic diffeomorphisms krol2004b . Thus the use of exotic $`R^4`$’s would be just respecting the philosophy behind general renormalization technology.
To be more specific, let us abbreviate the divergent counterterm in 10 by $`A(g_0,a_n,ϵ)`$. Taking the limit of $`ϵ0`$ we can reexpress it by taking infinitesimal small $`ϵ`$’s but the expression $`A(g_0,a_n,ϵ)`$, contains the symbol referring to the nonstandard numbers, $`ϵ`$, and can take nonstandard (infinite big) values: $`{}_{}{}^{}A(g_0,a_n,{}_{}{}^{}ϵ)`$. To recover standard expression for $`A`$, one can take standard part operation, or $`{}_{}{}^{}A`$ can be embedded into the weaker logical environment, namely one deriving from the smooth Basel topos. In this intuitionistic logic infinite big numbers can be considered as smooth finite. This expresses the fact that in the Basel topos the true real numbers contain nonstandard numbers as well. The classical projection, which is not unique, gives standard expressions or Robinsonian nonstandard ones. Remembering that we are in dimension four, we can make use of the MTSD exotic smooth structures. The structures in question can be generated by the dynamically changing language $`L(\mathrm{SET},)`$ (see Appendix B). Clearly, MTSD smooth exotic structures on $`^4`$ are invariant with respect to the changes given by the dynamical language. Hence, divergencies as in the countertem 10 should be the building ingredient of the exotic structure. However, how to describe the structure in terms of global differential calculus is not known at present.
The structures need the divergent terms as their building blocks and which are described as divergent in a specific, model–theoretic limit of the structures. Furthermore, the structures can compensate for the divergencies, hence one need not to subtract the terms on the 4–dimensional field theory side by hand. Instead, the theory should be formulated on exotic rather than standard 4-spaces.
Now we can see that exotic MTSD smooth $`R^4`$’s localized in the asymptotic $`AdS`$ boundary can generate the counterterms which, in turn, can cause that the bulk 5-dimensional gravitational action becomes finite. In addition to this the 5–dimensional gravitational divergencies are translated into exotic structures on 4-space, which shows the specific connection of exotica in 4–dimension and 5–dimensional gravity; thus, not only the 4–dimensional sources of gravity are to be considered as gravitational impact of 4–exotica. Besides, these are 5–dimensional gravitational divergencies. The MTSD 4-structures seem to be well suited for the renormalization questions.
However, more explicit calculations would require analytical tools which at the present time are out of our reach. This is the task for forthcoming studies which possibly will improve the situation. An important step towards grasping the Yang-Mills theories on some exotic $`R^4`$’s and towards description of exotic metrics on these, has been recently made by Kato kato2004 <sup>3</sup><sup>3</sup>3I would like thank professor Robert Gompf for letting me know about this paper. With the help of these metrics one could attempt to understand what is the role of exotic smooth $`R^4`$’s for the susy breaking pattern in the AdS/CFT correspondence. As a consequence, more realistic 4-dimensional YM theories can emerge as dual to the sustrings in 5 dimensions. More detailed studies of that case, however, have to wait untill exotic smooth invariants will be generated via model-theoretic methods. Although, some remarks concerning the modification of YM and sustrings sides of the Maldacena cortrespondence by exotic smooth $`R^4`$,s, one can find in krol2004 .
## Appendix A Relativity of formality and informality in mathematics
If we deal with some objects and try to interpret them as objects in some category, we would like to be able to deal with sets of the objects or sets of sets and so on. Simply, the objects should be freely arranged in sets with respect to some (meta) ZFC. These kinds of requirements are basic ingredients of what one may consider as a classical meta-environment for seeing the objects classically. So, being classical on the metalevel would mean being organized according to a set-like way. This would mean that the outer view of the category of the objects in question is regarded as classical provided it resembles the category SET.
To be more specific let us formulate the following working definitions:
* One says that the metalevel is organized classically if it is formally modelled by SET.
* Formally modelled by SET means the objects can be freely arranged into objects satisfying SET axioms.
* Freely arranged in sets can be understood in terms of set–based diagrams and its all finite limits or colimits which should exist, with the requirement that this outside SET structure agrees with the point structure of the objects while seen from the outside.
* We say a formal language is classical at the level of symbols if the symbols of the language are formally modelled by SET. The symbols of the language are interpreted in SET rather than in a general not two–valued and not classical topos. Thus, every symbol may have its name in some formal presentation with the use of classical sets similarly as the symbols are considered as separated and global in the preformal outside presentation.
* Similarly, we can say that a formal language is intuitionistic at the level of symbols provided the symbols of the language are internally interpreted in a topos. This is precisely the case where language or theory are interpreted in a topos moerdijkMcL1992 .
* Resemblance of some category with SET is to be understood as possesing almost all or all categorical properties of SET by some other category. The specific list of the categorical properties should be specified. Isomorphism between SET and some other category is the example of the resemblance of categories.
It is possible that a category resembles SET but does not respect the set–like structure of its objects. There are some toposes which have the categorical properties of SET, but the point–like structure of their objects is not preserved by morphisms of the toposes; the morphisms define elements and subobjects in the sense of the category but not in the sense of SET.
Let us consider, as the example, some non classical topos with a natural number object (NNO). Objects of the topos can be organized twofold:
* classically, as the objects ”from the outside” usually have classical set–theoretical structures, the one which agrees with classical meaning of element of a set and the description of them respects these structures.
* nonclassically or internally, according to the internal, not SET–like and non classical, topos structure. The topos structure enforces the understanding of being an element (member of a set) or being subobject (subset) etc.. In general, any topos generates internal set theory moerdijkMcL1992 ;
Even though the inherent nonclassicality derives from the not–SET and non classical structure of the topos exists, there is always the meta–level, such that the language of the theory of an elementary topos is the first order one. At this level one perfectly recognizes the symbols of the language as being arranged according to SET–like structures. The language is classical at the level of symbols. Moreover, at this level one can see the objects of the topos from the outside as having classical set-theoretical structures.
The basic intuition which is of our concern here is that in any kind of formal presentation of a theory there always exists a sufficiently high level of the metatheory which would be organized classically and SET–like. That means that the formal language of a metatheory would be classical at the level of symbols.
Although the intuition behind this statement supports its obviousness, we attempt to argue against it in the following sections.
Let $`L`$ be some formal language, of the first or the higher order, possibly with sorts $``$ or $``$ for real and natural numbers, and $`T_M`$ be a theory in the language which describes the manifold $`M`$.
To have formally presented all the manifolds from the category $``$ we have to associate to any manifold $`M`$ some formal language $`L_M`$ and the respective theory $`T_M`$ is formulated in this language.
If $`M`$ is considered as to be the universum of the model of $`T_M`$ we try to add nonlogical constant symbols to the language $`L_M`$ such that they exactly correspond to points of the manifold $`M`$. The intended interpretations of the constants in the model $`M`$ are just the points the constants correspond to. In this way we have what is refered to as a simple extension of the language.
We say that a formal language $`L`$ corresponds literally to some open subset $`U`$ of some manifold $`M`$ if the language has in the set of its constant nonlogical symbols the ones corresponding uniquelly to the all points from $`U`$. Then we use a symbol $`L_U`$ for such a language.
To see all the languages $`L_U`$ as sets of symbols we need some meta ZFC or PA to refer to, and $``$ as a sort is understood classically with respect to this ZFC as well. Thus, every real from $``$ may have a global name in the complex.
It is assumed that the reference to such a ZFC or PA is always possible, though in an informal way, and that this is the absolute feature of a sufficiently rich metatheory.
The family of sets of symbols $`\{|L_U|\}`$ determines the family of formal languages $`\{L_U\}`$ associated to the category $``$ and which are simple extensions of the language $`L`$.
###### Lemma 1
The family $`\{|L_U|\}`$ of sets of constant symbols corresponding to the formal languages $`\{L_U\}`$ is organized exactly as $`Sh()`$ that is, the sets of symbols behave as sheaves with respect to morphisms of the category $``$.
This is in fact expression of the well-known fact that the family of open subsets of a manifold form a Heyting algebra (etalè) rather than a Boolean algebra moerdijkMcL1992 and that the family of symbols of the languages should respect functorial morphisms in $``$ when taking finer open covers.
Now we know ReyesMoerdijk that there are full and faithful embeddings of $``$ in $`𝕃`$, in $`\mathrm{SET}^{𝕃^{\mathrm{op}}}`$ and in $`Sh(𝕃)`$. This means that
###### Lemma 2
The formal language which would describe literally open domains of manifolds from $``$ can be regarded as interpreted already in the topos $`Sh(𝕃)`$ (it cannot be classical at the level of symbols).
This is in fact reformulation of previous lemma and it states the fact that the families of sets of symbols are organized in the sheaves and $`Sh(𝕃)`$ allows for the SET–like external organizing the symbols. $`\mathrm{}`$
Taking extensions of the theory still allowing for the literal descriptions of the manifolds, as having constant symbols in the language literally corresponding to the open domains, gives the following
###### Corollary 1
In the case of the category of smooth manifolds and smooth maps between them any formal extension of the theory which would allow for literal descriptions of open domains of the objects of the category has to be intuitionistic.
This is because the formal and literal description in question has to be intuitionistic at the level of the language symbols.
The peculiarity of this corollary is in that the category in question is the one whose objects are well-defined classical objects, namely smooth real finite dimensional manifolds, and they certainly should have a classical language suitable for a formal description of them. But while we attempt to keep the classical outside view of all the manifolds, the functional incompleteness of the category of manifolds become important and in the case of literal formal descriptions the language cannot be classical any more. This is the theory–metatheory complex which allows to see the language symbols as sets.
If we instead allow less formal connection between sets of language symbols and objects in the category (as usually is the case) then the informality compensates nonclassical but formal character of the language symbols in question, and the symbols can be considered as classical sets. Moreover, the use of sorts $``$ or $``$, requires classical and set–like understanding of the sorts, if formally grasped, hence any element of the sorts should have its global name. Thus, extending the languages describing the manifolds toward literal ones is quite natural requirement while the sorts are in use.
###### Corollary 2
Any higher order extension of the theory, describing literally open domains of the manifolds from $``$, has to be performed internally in the topos of sheaves of the language symbols.
This is just the reformulation of Lemma 2. $`\mathrm{}`$
If, instead, one considers the formal theory of manifolds which is not literal in the sense above, the extensions of such a theory can be classical and are not necessarily performed internally in the topos. In fact, literal formal presentation of the objects of the category and the language classical at the level of symbols cannot go together.
If the objects of the category are considered formally and classically according to some meta ZFC then the languages describing them have to be well-defined in the sense that they have to have sets of symbols given classically. But in the case of topos–like modeling of the symbols of the languages we do not have the symbols as classical entities. Moreover, the topos can be described formally in the first order language classical at the level of symbols. Thus, the following holds:
The metalevel in the theory–metatheory complex has a formal intuitionistic description, extending the literal theory with some topos of sheaves of the symbols of the formal languages but this description is non-simultaneous with the formal descriptions of the objects of the category given by the languages.
Let us note that in the case when the metalevel is not a formal extension of the theory, the relation of the theory and the metatheory has to deal with some non formal ingredients. In this situation there is a possibility that the metalevel is presented formally as the formal topos and expressed in the language classical at the level of symbols. This presentation can be performed along with the formal presentations of the objects of the category. In this way we have some means allowing to distinguish between formal connection of the metalevel and the theory and a nonformal connection. Here is where informal ingredients still have to be presented in the complex.
Hence, we have the higher order extensions of the literal theory of manifolds from $``$ constructed internally in the topos, or else classical SET–like environment where manifolds are arranged classically as language symbols do.
In the first case, the metatheory is topos theory and, if formal, also this metatheory has its symbols organized classically. In the second case the metatheory is just ZFC or part of it. So, there is the classical ZFC meta-environment where classical manifolds are not within the reach of a literal and formal description. In the opposite case there is the intuitionistic meta-environment where manifolds can be reached formally and literally.
One can say that when informal ingredients are allowed, the manifolds are seen classically but, when formally and literally, they have to be presented internally in the topos. In fact, we claim that both situations are not independent as purely objective, well-separated pictures of the same unaffected reality. As was conjectured in krol2004 , there are objects in the category $``$ which require both of the pictures being entangled.
## Appendix B Intuitionistic versus classical metaenvironment. Dynamics of the changes
Let us, following MacLane and Moerdijk moerdijkMcL1992 , recall some main concepts connected with a language interpreted in a topos.
Given any topos $`𝒯`$ we can specify the language of the topos which would allow us to built internally set–like constructions and it could express the true statements about the topos. The language is called the Mitchell–Benabou language. All constants and variables as well all formulas are of some type. The types are considered as objects of the topos $`𝒯`$. For each type $`X`$ there are variables of this type and they are represented by the identity morphism $`1:XX`$.
A term $`t`$ of type $`X`$, possibly containing some free variables $`y,z,w`$ of types $`Y,Z,W`$, is represented by an arrow in $`𝒯`$
$`t:Y\times Z\times WX`$
We do not present here detailed inductive definition of terms and their interpretations in a topos, since we do not use them explicitly. Again, all details can be found in moerdijkMcL1992 .
Terms $`\varphi ,\psi ,\mathrm{}`$ of type $`\mathrm{\Omega }_𝒯`$ (subobjects classifier in the topos $`𝒯`$) are formulas of the language. We can connect formulas with various logical connectives and apply the quantifiers which yield formulas again of type $`\mathrm{\Omega }_𝒯`$. This is possible, since the logical connectives are simply Heyting algebra operations defined on $`\mathrm{\Omega }_𝒯`$, and can be composed with the morphisms ending on $`\mathrm{\Omega }`$ as is in the case of formulas.
Similarly, the definition of quantifiers is possible internally so that while the quantifiers are applied to the formulas it results in the morphisms composition again of type $`\mathrm{\Omega }_𝒯`$. Now, let us slightly modify the above $`Sh(𝕃)`$ category such that the Grothendieck topology $`J`$ on it is going to be changed into the one which is also subcanonical (all representable presheaves from the Yoneda embedding are sheaves) and NNO is the one which contains infinite nonstandard naturals as well. This kind of modification is possible and relies only on the changes in the topology leaving the objects in $`𝕃`$ unaffected. The resulting category of sheaves in the changed topology on $`𝕃`$ is called Basel topos ($``$) and it was described in details in Moerdijk and Reyes ReyesMoerdijk .
Again, the sheaves in $``$ can be considered as sheaves of language symbols, but with respect to new covering families deriving from the changed topology. The classes of sets of symbols which give descriptions of objects of $`𝕃`$ are organized as sheaves on the category $`𝕃`$. This results in the structure of Basel topos $``$.
The change in the topology allows us to take as covers for the topology the ones which generate the Grothendieck topology as before, but with the addition that all projections along all nontrivial loci generate the covers as well ReyesMoerdijk .
Now, we attempt to be more explicit in demonstrating the role played by interpreted language in the topos while the special kind of classical objects is considered. The basic observation is that the object of natural numbers in $``$ contains infinite ones. The classical counterpart of this is the existence of nonstandard models of PA (in the sense of Robinson) which also contain the infinite big natural numbers. Furthermore, the object of smooth real numbers in $``$ contains (intuitionistically) indempotent reals $`d^2=0`$ and infinite big and infinite small ones. Classically indempotents reals do not exist in any model of real numbers but again, there are Robinsonnian models of real numbers which are classical and which contain nonstandard big and small numbers. Thus, one can consider the nonstandard models of real numbers as a classical counterpart of intuitionistic object of smooth real numbers. One can say that the object of smooth NN and smooth real numbers from $``$ are projected on classical domains (or projected classically), namely on the $`{}_{}{}^{}𝐍`$ and $`{}_{}{}^{}𝐑`$ respectively.
Let us introduce some working definitions.
* An object exists classically provided its existence is verified by some classical theories.
* Classical theories are those which are based on classical logic but they are not necessarily axiomatized or even purely formal. They should be considered as the complex theory – metatheory which gives mathematical results (theorems).
* A formal language, $`L(\mathrm{SET},)`$, is called to be dynamically changing from $`\mathrm{SET}`$ to $``$ if among its symbols are those from $`\mathrm{SET}`$ and others internally interpreted in $``$.
* We say that an object which exists classically is described generically by a dynamical language $`L(\mathrm{SET},)`$ if the object has to refer both to the numbers from $``$ ($``$) and from $`{}_{}{}^{}𝐍`$ ($`{}_{}{}^{}𝐑`$) by the symbols of the language, but $`{}_{}{}^{}𝐍`$, $`{}_{}{}^{}𝐑`$ are considered as projected classically from $``$.
###### Lemma 3
If the meta-language is considered as changing dynamically from the constant SET-like into the one interpreted internally in Basel topos, and if there exists a classical object described generically by the varying language, then there exist two models of PA, one standard and one nonstandard given by the ultrafilter construction, which are not formally distinguished from one another. This means that there does not exist any higher order property of standard natural numbers, which would not be valid in the nonstandard model.
This is possible since the higher order properties are defined with respect to SET and compared with those defined with respect to the internal language of $``$. If in SET, standard model, $``$, of PA appears. If internally in $``$, then classical projection of $`smooth`$-$`N`$ gives $`{}_{}{}^{}𝐍`$. $`{}_{}{}^{}𝐍`$ is given classically by the ultrafilter construction.
Based on Lemma 3, we have
Under assumptions of Lemma 3, the standard model of real field and some nonstandard one given by the ultrafilter construction are indistinguishable. This means that there does not exist any higher order property of standard real numbers, which would not be valid in the nonstandard model. This indistinguishability has a formal meaning for some smooth manifolds from $``$, namely for exotic $`R^4`$’s.
Similarly, one cannot distinguish the standard real numbers and nonstandard ones, by any higher order property, provided the quantification is performed internally i.e. over internal subsets of $`{}_{}{}^{}𝐑`$ kock1974.
Formal means needed to distinguish the models are already internally interpreted in $``$. On the other hand, informal ingredients have to be always present in the complex theory–metatheory, and these sometimes cause lack of formal control over the meta-environment. The possible meaning of the indistinguishability $`^4`$ and $`{}_{}{}^{}𝐑_{}^{\mathrm{𝟒}}`$ for some (classical) manifolds is conjectured krol2004 ; krol2004b . The smooth structures of the manifolds can survive the shift. Exotic $`R^4`$’s are certainly classically existing entities, since the proof of the existence of these was obtained classically.
Equivalently, in terms of invariant mathematics, i.e. valid in all toposes with NNO, one can say that formal differences between local presentations with respect to the specific frames-toposes are neglected. Only properties valid in all toposes with NNO are expressible invariantly. In particular, in $``$ NNO is $`s`$-$`N`$ and in SET this is standard $``$ and the invariance requires that no difference between them is expressed. This indistinguishability, however, can have in turn, formal back-reaction. Thus, the shift from SET to $``$ has a price in mathematics; it is the emergence of exotic $`R^4`$’s. And this is like the emergence of gravitational effects in the covariant, tensorial, language of GR.
One can also say that there are some smooth manifolds requiring both classical and intuitionistic meta-levels if formal.
The indistinguishability of the models of PA or real numbers is refered to as $`main`$ $`hypothesis`$ in what follows.
To be more specific, let the types of the language symbols be sheaves in $`Sh(𝕃)`$ or $``$. It means that the variables and all terms are interpreted by suitable morphisms in the topos $``$. All formulas are as usual of type $`\mathrm{\Omega }_{}`$ (classifying object of $``$).
The entire interpretation is forced by
* Taking constant symbols as being of the type of some sheaves in the topos.
* Respecting the SET–like, classical grouping of the symbols of the language i.e. respecting the inverse limits etc. as in the categorical presentation.
Thus, the variables are also of the $``$–sheaves types. Subsequently, one has to interpret the whole language in the topos. Let us note that a topos can serve as a model for any higher order language.
In this way we can have all set–like classes as $`\{x:\varphi (x)\}`$ already in the topos, and the validity of formal statements in the topos can be translated into the outside one simply by avoiding the excluded midle law and the axiom of choice in the argumentation performed in the internal language. In fact, we have a language of the invariant mathematics. Being equipped with such an internal language in $``$ we have to use it in the context of smooth manifolds which are formally and literally approached.
The special objects can exist which are able to survive as classical objects the shift in the descriptions given by constant classical language and the language already interpreted in the topos as above. In the case of smooth manifolds this would result in the situation where both descriptions should be connected by some diffeomorphisms if the object survives as the classical object. Since the diffeomorphism cannot be standard, this may serve as an approach toward the understanding of the so called exotic diffeomorphisms in the case of exotic smooth $`R^4`$’s as in krol2004b .
###### Acknowledgements.
I would like to thank professor Robert Gompf for the correspondence and explaining to me many points regarding exotic $`R^4`$’s. I owe special thanks to the organizers of the IPM String School and Workshop for their hospitality and giving to me the chance to participate in the School.
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# On the Job Training
## 1 Introduction
This paper proposes a new way of modelling machine learning (ML), inspired by stream-based active learning but more accurately reflecting strategies humans employ when working with trainers and more directly addressing the needs of users reluctant to risk effort integrating ML algorithms. Like active learning it aims to reduce labeling overhead, but unlike existing models, it removes the separation between training and testing, leading to systems which can quickly adapt to their tasks as well as receive ongoing training throughout their useful lives. We call our approach “on the job training” (OJT).
Traditional paradigms assume that an administrator prepares training and test sets with distributions similar to that of the real problem, assesses performance of the agent, then deploys a fixed instance of the classifier once it reaches a satisfactory level of performance. However, this train/test/use paradigm does not adequately model the needs of systems which must begin performing immediately, or living systems whose datasets change and for which additional classified data becomes available over time. The OJT framework explicitly addresses this by formally defining a concurrent train/use process along with performance metrics which indicate both immediate and long term performance. Traditional classifications distinguish between active and passive systems, and within active learning, between pool-based and stream-based designs. We expand this space to include synchronous and asynchronous relationships between agent and trainer, the presence or absence of intervention from the trainer, and temporally evolving distributions of interest.
This is a framework paper exploring one particular niche in a spectrum of learning paradigms, which we claim reflects the needs of many real-world practitioners. Like stream-based active learning, our model lets the agent decide how confident it feels in its classification ability, allowing it to take over for the trainer when easy problems arise or to ask for help when a difficult query comes along, maximizing the value of the trainer’s time. Our performance metrics are unique in measuring how rapidly the agent gains the ability to relieve the trainer from tedious tasks without neglecting long-term accuracy. We also assume that the trainer is willing to work interactively with the agent, giving unsolicited advice or allowing the agent to specify what points it would most like to have classified, as a pool-based active learner would. These points will often identify interesting or problematic areas for the trainer to research.
Consider how people tend to evaluate software. Users almost never invest a large amount of initial effort reading a program’s manual, prefering instead to start working immediately to see if the software does what they need. While they have some initial patience, they expect to see results quickly, and if the package isn’t working when their patience runs out, they’ll discard it, regardless of how it might have done in the long run. Conversely, users also expect software that they use over time to remember their changing preferences, and get frustrated with software that continually requires overrides to its initial settings.
Our metrics are designed to address these human requirements, emphasizing the need to solve initial problems correctly with human aid, minimizing up-front training demands, yet still requiring satisfactory results in the long run – in short, suitable for on-going, real-time interaction with humans.
Consider also the problem of constructing security policies. Low security installations tend to use overly broad policies rather than expend the resources required to handle the numerous exceptions and permission grants which occur over time. Administrators will be unwilling to spend time with ML assistants unless they can adapt to their task very rapidly, responding to correction, asking for help only for boundary conditions and adapting to later policy changes. High security installations have to balance the need for tight control with the danger of overly complex policies which will have subtle errors. Their administrators are unwilling to allow an agent to make security decisions, but can still benefit from an agent which can point out parts of the access space which seem most unusual. In both cases, users want an agent which can quickly give useful results, yet perform in the long term as conditions change.
There is little work directly related to this paper. Active learning traditionally focuses on different *learning techniques* (such as the relative advantages of certainty based methods compared to committee based methods ), as opposed to exploring new *performance requirements*. Closest in spirit to OJT are stream-based active learners , but these have generally been evaluated in terms of overall classifier accuracy, rather than with respect to their ability to address the immediate human needs we have illustrated.
In addition to defining our class of agents and corresponding performance metrics, we will specify how to build an ideal but impractical OJT learner, and consider an obvious approximation to that algorithm based on KNN. We conclude with promising performance results and discuss future research.
## 2 System Specification
Define $`S`$, a sample set of zero or more unclassified points which may be used as questions by the selective sampling algorithm (systems which use membership queries instead of selective sampling may leave $`S`$ unspecified). $`S`$ is provided to the agent at the start of the execution phase. Define $`C(q)`$, a function which returns the true classification of any point $`q`$. Let $`G_t(q)`$ denote the agent’s “guess” at the classification of any point $`q`$ at a particular “time” $`t`$. Let $`V=\mathrm{}`$, the set representing the queries which have arrived. Let $`t`$ be the loop iteration counter. Define $`Perf(t,q,Q,f,G_t(q),C(q))`$, an (optional) performance metric which keeps track of the agent’s classifier throughout the execution phase.
An OJT agent executes a series of five basic steps, which may be arranged differently depending on the system design. The steps are 1) Query. A query $`q`$ is sent to the agent. Set $`V=Vq`$. $`q`$ can be null for asynchronous agents, meaning that no query needs to be answered at this time. 2) Question. The agent outputs a question $`Q`$, an unclassified point it wishes the trainer to classify. In selective sampling implementations, $`QSV`$. $`Q=\mathrm{}`$ implies the agent chooses not to ask a question. 3) Fact. The trainer sends the agent $`f=Q,C(Q)`$, or $`f=\mathrm{}`$ if $`Q=\mathrm{}`$. $`f`$ can be null for asynchronous agents, meaning that no fact is forthcoming at this time. Real systems might also allow the trainer to set Q, providing an unsolicited fact. 4) Answer. The agent outputs $`G_t(q)`$. 5) Assess. Call $`Perf(t,q,Q,f,G_t(q),C(q))`$ to assess the agent’s performance.
Here, we make our first distinctions between agents. A *synchronous* system requires a query to be provided at each loop iteration, and that the trainer always answers questions asked by the agent during the iteration in which they’re asked. Note that agents are not required to ask a question during a particular loop iteration. Real world problems which benefit from on the job training are likely to have an *asynchronous* nature, however, in which queries may arrive at any time, and in which trainers may not always be available or willing to answer questions from the agent. Furthermore, the facts provided by the trainer are not required to correspond to the questions asked by the agent.
The order of operations in the loop is critical, and motivates the distinction between systems with and without *intervention*. If the answer phase is placed after the fact phase, the agent has no opportunity to choose its question so as to aid in answering the query at hand, and is said to be a system without intervention. If the query set $`V`$ follows a different distribution than the sample set $`S`$, then the agent can still use $`V`$ to specialize on part of the problem space. But if the distributions of $`V`$ and $`S`$ are the same, then the agent in a system without intervention has no opportunity to improve its accuracy for an unanswered query.
An agent with intervention therefore executes the steps in the following order: Query, Question, Fact, Answer, Assess. An agent without intervention executes the steps in this order: Query, Answer, Question, Fact, Assess. Synchronous and asynchronous systems are distinguished by whether or not null queries and facts are permitted.
## 3 Performance Metrics
Since OJT aims to encourage creation of both agents which respond quickly to new tasks and agents which are used and trained over long periods of time, we propose two new performance metrics. The first fixes all but two variables, giving the algorithm a fixed budget of questions to ask. The second metric uses an arbitrary utility function, allowing question and wrong answer values to be externally provided. A crucial point is that these metrics are made available to the agent, which allows it to optimize with respect to the performance measure it is being evaluated by.
Both metrics calculate agent accuracy cumulatively, as the queries are answered. Because of the lack of a formal testing phase, we calculate accuracy using the answers provided at each step. The focus on accuracy from the start requires a finite limit on the number of rounds considered if there is also to be a finite limit to the number of questions. Otherwise, as we will show later in this section, cumulative accuracy becomes indistinguishable from traditional accuracy measurements.
For the first metric, this requires the user to decide how important the agent’s learning rate is as a function of the amount of training he is willing to provide. This is natural decision – a user trying to solve a problem in one hour wants an algorithm which will very quickly be of some assistance, while a user preparing an agent for high volume, long term service will be willing to offer more training to an algorithm in exchange for high overall accuracy.
Budget Metric. Our first metric mirrors the traditional active learning model in which an algorithm chooses a preset number of points for classification. Since this metric’s utility lies in the simplicity of choosing only question and query limits, we specify it with respect to a synchronous system (but it is readily adapted to other variants). The metric computes the agent’s cumulative error over the course of execution. Agents are expected to use all available questions; no reward is given for unused questions at termination.
Define $`k_q`$ to be the query limit and $`k_Q`$ to be the question budget. Let $`c`$ be the counter for incorrectly classified queries, and $`Q_{rem}=k_Q`$ be the number of questions remaining. The budget metric is then defined as the cumulative error rate, or $`c/k_q`$ (when $`t=k_q`$).
This simple measure would be no different from traditional error calculations except that $`c`$ is accumulated after each query is answered, in the Assess phase. In section 4 we give two theorems which demonstrate the significance of this seemingly small difference. Note that, since this metric considers the *rate* at which each learner learns, our metric is quite useful for evaluating active learners, since their strength is generally considered to be in their fast convergence.
Utility metric. In the utility metric, the Assess phase calculates $`c`$, the cumulative cost as $`\sigma C_Q(Q_t)+C_w(G_t(q))`$, where $`C_Q`$ and $`C_w`$ denote the cost of asking a particular question or being wrong about a particular query. If possible, these cost functions should be made available to the agent. At the risk of being too general, this metric illustrates that agents should take advantage of questions that are easy for a trainer to answer (or times when the trainer is available to work with the agent) and should be willing to guess when a wrong answer will not be expensive.
## 4 Intervention and $`k_q`$
In this section, we show how two of the seemingly minor elements of our system and metric definitions have a dramatic impact on agent strategy and performance. First, we show that systems without intervention have advantages over traditional active learners only when the distribution of queries over time reveals information about the task at hand. Second, we emphasize that OJT’s distinctions from traditional active learning fade once questions are no longer being answered, and that consequently, the choice of $`k_q`$ in our proposed metrics is critical in making useful assessments of an OJT system.
###### Theorem 4.1
Assuming queries are chosen uniformly at random, with replacement, from a set $`T`$ whose distribution is the same as a sufficiently large unlabeled sample set $`S`$, an agent in a synchronous OJT system without intervention has no statistical advantage over a traditional active learning agent with respect to the budget metric.
Proof: Consider the first iteration for an agent $`A`$. $`A`$ is forced to answer a query $`q`$ before receiving the classification for any point. Once it provides its (random) guess, it is allowed to choose any point as its question. Since the trainer returns a fact only after the agent has output its guess, the fact clearly cannot influence that guess. But since $`q`$ was chosen randomly from $`T`$, and $`T`$ has the same distribution as $`S`$, it provides no information, in a statistical sense, about any future query. Consequently the agent can discard $`q`$ after guessing its classification. But this, again, happens before the question phase, so $`q`$ provides no useful information about what question to ask. By induction, we see that the agent has the same dilemma for all future iterations of the loop. The agent can therefore equivalently output its classifier $`G_t`$ in each iteration of the loop before the query arrives, then ignore the query once it does, since the query provides no utility to the agent for current or future rounds. Thus $`A`$ can equivalently operate as either an OJT agent or a traditional active learner with respect to the budget metric. $`\mathrm{}`$
In the event that $`T`$ and $`S`$ have different distributions, OJT agents have an advantage over active learners even in systems without intervention. But also note that this proof doesn’t apply to systems with intervention, in which OJT can outperform baseline even for problems which are traditionally considered unlearnable.
### 4.1 The importance of $`k_q`$
While OJT systems should generally produce agents which have good overall accuracy, their emphasis is on learning the query set at hand. While this seems natural for many applications, it is useful to note that if all you want is a system to train once and then deploy, OJT has little to offer over traditional active learning, as this theorem demonstrates.
###### Theorem 4.2
Define an optimal agent as one whose classifier accuracy monotonically increases as it learns queries and facts, and which achieves cumulative error less than or equal to the cumulative error of any other practical classifier. Then given a finite $`k_Q`$ and an infinite query set $`T`$, the difference in cumulative error $`ϵ`$ between optimal agent A in an OJT system with and intervention and optimal agent B in an OJT system without intervention approaches zero as $`k_q`$ approaches infinity.
Proof: Assume that an optimal agent $`A`$ in an OJT system with intervention can achieve a cumulative error of $`e_A`$ by asking no more than $`k_Q`$ questions over $`k_q`$ queries from $`T`$. An optimal agent $`B`$ in a system without intervention can receive precisely the same information available to $`A`$, but only at a later time, so its cumulative error $`e_B`$ can never be lower than $`e_A`$. Let $`ϵ=e_Be_A`$. Assume that $`ϵ`$ is maximized, so that $`A`$ answers $`k_Q`$ queries correctly which $`B`$ misclassifies as a result of not receiving facts until later in its loop. Then for a given $`k_q>k_Q`$, $`ϵ=k_Q/k_q`$, which approaches 0 as $`k_q`$ approaches infinity. $`\mathrm{}`$
Despite our questionable definition of an “optimal agent,” it should be clear that the advantages of intervention fade as one considers cumulative accuracy past the point where the question budget has been exhausted. This should not be taken to mean that OJT is no different from other types of ML; rather, it should emphasize that OJT’s strengths lie in adaptability to temporally evolving conditions and ability to minimize load on trainers.
## 5 Implementation
Here we show how to create both an ideal, but impractical OJT agent as well as a practical agent which generally outperforms its active learning counterpart. Both implementations assume a budget metric and synchronous systems with intervention.
Modeling OJT agent strategies is easier if the agent can predict how accurate it will become if given the classification of a particular unlabeled point. The implementations we propose all assume that the underlying classifier has the ability to recursively predict how it will classify future points, and how confident it will be in that classification. Let such an algorithm provide three functions: $`Unc(p)`$ returns the probability that p will be misclassified. $`Add(p)`$ assumes that $`p`$’s label is known when calculating $`Unc()`$ ($`Add`$ may be called more than once to assume multiple points are known). $`Remove(p)`$ means that the algorithm should no longer assume that $`p`$’s label is known.
### 5.1 Ideal OJT
Given an agent with perfect foresight and unlimited computing power, and some simplifying assumptions about the system, it is straightforward to construct an ideal system in the sense that it minimizes cumulative error with respect to the budget metric. We begin by assuming the entire test set $`T`$ and remaining test set $`R`$ are known, and that queries are chosen from $`R`$ uniformly at random. Let $`q`$ be the query just selected in the query phase (which is no longer in $`R`$), and $`Q_{rem}`$ be the number of questions remaining in the question budget. Let $`Permute(T,f)`$ calculate all permutations of the elements of $`T`$, calling $`f`$ on each permutation; let $`Average(S)`$ returns the average of the elements in set $`S`$. Then the ideal OJT algorithm for a selective sampling synchronous OJT system with intervention is as follows. Define a function $`ExpectedWrong`$, which returns the number of queries it expects to misclassify by the time $`R`$ is exhausted (which is minimal on average given our assumptions and its choice of $`Q`$). Assume $`ExpectedWrong`$ is called during the question phase of system execution, and that the value $`Q`$ is the question it selects (we omit the corresponding algorithm which optimizes the utility metric rather than the budget metric, which is easily derived). Letting the function AvgPenalty$`(R,Q_{rem})`$ return $`Average(Permute(R,ExpectedWrong(R[0],RR[0],Q_{rem})))`$:
ExpectedWrong$`(q,R,Q_{rem})`$:
If $`q==R==\mathrm{}`$: Set $`Q=\mathrm{}`$; Return $`0`$.
// Start assuming we ask no question
Let $`Q=QBest=\mathrm{}`$
Let $`MinPenalty=Unc(q)+AvgPenalty(R,Q_{rem})`$
// Then consider all possible questions we might ask
If $`Q_{rem}==0`$, Return $`MinPenalty`$
Loop for all $`Q(ST)`$:
Set $`Penalty=Unc(q)+AvgPenalty(R,Q_{rem}1)`$
If$`(MinPenalty>Penalty)`$
Set $`MinPenalty=Penalty`$
Set $`QBest=Q`$
Set $`Q=QBest`$; Return $`MinPenalty`$
Note that since we assume questions are chosen randomly from $`R`$, we can average the penalties over all possible futures, choosing the question in this round which leads to an optimal outcome on average, since all sequences are equally likely.
Note that this algorithm is entirely intractable; $`Permute`$ calls $`ExpectedWrong`$ for every possible future, and $`ExpectedWrong`$ is itself a recursive function. But it does suggest a general form which a practical algorithm might take if it can approximate or efficiently calculate the expected penalty across reasonable futures.
### 5.2 A Tractable Approximation
The KNN classifier described in provides an excellent starting point for implementing OJT classifiers. Their classifier implements active learning using selective sampling, and outperforms similar passive learners as well as several active learning approaches adapted from other sources; it also makes it simple to construct the $`Add`$, $`Remove`$ and $`Unc`$ functions required by our implementation. In fact, conditional uncertainty with lookahead is the basis for their utility metric.
Our first approximation of the ideal algorithm effectively answers the question, “if I have only one question to ask, and I ask it now, what should it be?” It approximates this implementation of $`ExpectedWrong`$ by substituting $`SV`$ as an approximation to $`R`$. Using $`Unc(q|Q)`$ as a shorthand form of $`Add(Q);Unc(q);Remove(Q)`$, the algorithm is:
$$\underset{QSV}{\mathrm{min}}Unc(q|Q)+\frac{q_{rem}Q_{rem}}{|R|}\underset{pR}{}Unc(p|Q)$$
The summation of uncertainties divided by $`|R|`$ gives the expected overall classifier error given the true classification of the $`Q`$ under consideration. Multiplying by $`q_{rem}`$ would give the expected number of misclassifications given no further questions, but first we subtract $`Q_{rem}`$, since at least that many additional queries can be classified correctly just by passing them along as questions. We then add the penalty of misclassifying the current query explicitly, since a penalty will certainly be incurred if it is misclassified, as opposed to the possible future queries whose uncertainty is amortized. Most obviously missing from this implementation is the option of asking no question at all in a given round, since we have not yet found a reliable, efficient way to predict when this behavior will be worthwhile.
### 5.3 Performance
We compared our classifier to two others. The first uses the simple strategy of using the first $`k_Q`$ queries as questions as they arrive. This ensures 0% cumulative error for the first $`k_Q`$ rounds and performs as a passive learner would for the remaining queries. The second strategy is a simple adaptation of the active learning algorithm proposed in . It adds queries to its unlabeled set as they come in, but otherwise behaves normally. In all cases, we used the budget metric and an overall accuracy metric with $`k_Q=20`$ for evaluation purposes, since the results in used approximately 20 training examples for each classifier they tested and because the budget metric is easier to directly compare with overall accuracy metrics. Note that we were unable to reproduce the exact results in .
Figure 1 shows our results on the ionosphere dataset from UCI . The active learning algorithm’s representation in the overall metric is not entirely fair, since the overall metric uses the same queries sent to the agents during execution. An entirely separate test set would give a better picture of the true generalized accuracy of each classifier, although our metric is a good halfway point between such a general metric and our cumulative metric. But as our second theorem points out, as $`k_q`$ increases, such advantages disappear anyway. This can be seen on the right graph, where the cumulative metric more closely matches the overall metric and the gap between the active and OJT learners decreases slightly. We expect that gap to disappear entirely given a large enough $`k_q`$. Correcting for this inequity analytically, we see that the results are also consistent with what we expect. Given the active learner’s 17.6% overall accuracy, we would expect it to be wrong on about 3.5 of the first 20 queries, while the OJT agent might get all 20 correct. The test set used had 150 elements, so we expect the OJT learner’s results to be about 2.3% higher than they would be on a completely independent test set would indicate. Since this is almost exactly the gap between the active and OJT learners, we suspect such a test would show that their long-term accuracies are almost identical.
Figure 2 shows how the OJT agent quickly adapts to the query set, achieving a low cumulative error rates by the time the question budget is exhausted. Finally, Table 1 lists our results from two other datasets. The Image Segmentation dataset produced results comparable to the Ionosphere set, although the active learner edged out the OJT agent even with the handicap discussed earlier. The Pima Indians database was surprising, however, in that the “always ask” strategy outperformed both other agents in both cumulative and overall accuracy; investigating this result is left for future research.
## 6 Conclusions
This is a wide-open research area. Our practical implementation provides an initial attempt at maximizing the goals of OJT, and can be vastly improved. Asynchronous systems combined with creative cost functions in the utility metric provide a host of possible directions for further research, including consideration of costs which vary over time, and burstiness of fact input when trainers are only intermittently available in asynchronous systems.
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# Halo Mass Profiles and Low Surface Brightness galaxies rotation curves
## 1 Introduction
In recent years the inner structure of dark matter haloes has been the topic of some discussion. Observations suggest that the dark matter distribution in disk galaxies has a roughly constant density core, with a typical size of order a few kpc. A recent analysis suggests a shallow power-law mass density distribution $`\rho (r)r^\alpha `$ with $`\alpha 0.2\pm 0.2`$ (de Blok et al., 2003). On the other hand, numerical simulations based on the ($`\mathrm{\Lambda }`$)Cold Dark Matter (CDM) paradigm suggest a very steep inner mass density distribution, a so-called “cusp” (Navarro, Frenk & White, 1996, 1997) \[NFW\]. The most recent simulations suggest slopes with $`1.5\alpha 1`$ (Fukushige et al., 2004; Klypin et al., 2001; Ghigna et al., 2000; Moore et al., 1998, 1999; Jing et al., 1995). There is strong consensus that this slope cannot be as shallow as the observations suggest.
The rotation curves of Low Surface Brightness (LSB) galaxies are considered particularly clean tests of the CDM paradigm at galaxy scales. LSB galaxies are dominated by dark matter (de Blok & McGaugh, 1997), and their dynamics should thus give a more or less direct map of the dark matter distribution. LSB galaxy rotation curves, together with the missing dwarf or substructure problem (e.g. Moore et al. 1999), currently form the most serious challenges for the CDM model. There has been much debate on this topic in recent years. This is however not the place to review this debate, and the reader is referred to the summaries by e.g. Binney (2004); de Blok et al. (2003); Swaters et al. (2003) and de Blok (2004) and references therein.
On the theoretical side the “universal mass density profile” (with its characteristic $`\alpha =1`$ cusp) as proposed by Navarro, Frenk & White (1996, 1997), has long been a corner-stone of CDM (but see Moore et al. 1998, 1999). Recent simulations presented in Navarro et al. (2004) and Hayashi et al. (2004) \[hereafter H04\] show that simulated CDM haloes may actually contain a range of slopes in their inner regions (though the slopes remain steep). The shallowest inner slope measured is $`\alpha 1`$, with an average over the simulated haloes of $`\alpha 1.2`$ (at the innermost reliably resolved radius for galaxy mass haloes). With this scatter in mind, H04 argue that simple fitting functions (such as “the” NFW halo) do not capture the full variety and diversity in shapes and slopes, and suggest that this may partly explain the discrepancy between observations and simulations. They claim that the inner slope is difficult to constrain, both observationally as well as theoretically and point out that rotation curve constraints are strongest where numerical simulations are least reliable. This, as well as the potential for systematic effects in the observations, has, at least according to H04, led to an “unwarranted emphasis” on the value of the inner slope, at the cost of evaluation of the data and simulations over their full radial extent.
H04 present a novel method to fit rotation curves derived from their simulations as well as LSB galaxy rotation curves from the literature. Their method is able to describe the variety present in data and simulations and from the distributions of respective fitting parameters H04 conclude that the observed curves *are* consistent with the CDM paradigm; a conclusion that contradicts a large body of observational work.
For their analysis, H04 make use of the rotation curves presented in McGaugh et al. (2001), de Blok & Bosma (2002) and Swaters et al. (2003), and derive circular velocity curves for haloes from their simulations. To quantify both sets of observed and simulated rotation curves they use a three-parameter fitting formula (see e.g. Courteau 1997):
$$V(r)=\frac{V_0}{\left(1+\left(\frac{r_t}{r}\right)^\gamma \right)^{1/\gamma }}.$$
(1)
Here $`V_0`$ is the asymptotic velocity of the flat part of the rotation curve, $`r_t`$ is a scale radius (the transition radius between the rising and flat part of the rotation curve) and the parameter $`\gamma `$ describes the abruptness of the turn-over between the rising and flat parts of the rotation curve. A higher value of $`\gamma `$ results in an abrupt transition, a low value in a more gradual turn-over. Note that this function was designed to fit observed rotation curves: it specifically assumes a solid-body rise in the inner parts and a flat rotation curve in the outer parts.
H04 fit the two sets of rotation curves with this function, constraining the fit parameters slightly by insisting that
1. $`0<\gamma 5`$,
2. $`r_t>0`$,
3. $`V_02V_{\mathrm{max}}`$,
with $`V_{\mathrm{max}}`$ the maximum observed velocity in the rotation curve. Fits with $`\gamma >5`$ \[rejected by condition (1)\] correspond with very abrupt transitions between the rising and flat parts of the curve. Condition (3) prevents some of the solid-body rotation curves where the asymptotic velocity is not well-constrained, from producing fits with extremely large $`V_0`$. Conditions (1)–(3) thus help removing unrealistic values for the fit parameters.
Figure 1 (after H04, their Fig. 9) shows the results presented in H04. The LSB rotation curves show a broad distribution in $`\gamma `$ with two pronounced peaks: one around $`\gamma 1`$ and a narrower one at $`\gamma 5`$. Some 70 per cent of the fits have $`\gamma 2`$. The $`\gamma `$ distribution of the simulated CDM haloes is markedly different. It is narrow and centered around $`\gamma 0.6`$ with a dispersion of $`0.4`$.
At first sight the two distributions seem quite incompatible, and one might conclude that the LSB rotation curves are not consistent with the CDM simulations. H04 argue however that for most of the LSB galaxy rotation curves the reduced $`\chi ^2`$ distribution is quite broad and shallow, so that by constraining the fit parameters slightly more, a fit can be found that is consistent with the simulations, at the cost of only a small increase in $`\chi _{\mathrm{red}}^2`$. They define these CDM-compatible constraints as:
1. $`0<\gamma 1`$
2. $`r_t>0`$
3. $`V_02V_{\mathrm{max}}`$
4. $`\left|\mathrm{log}\mathrm{\Delta }_{V/2}\mathrm{log}\mathrm{\Delta }_{V/2,\mathrm{CDM}}\right|0.7`$.
Here $`\mathrm{\Delta }_{V/2}`$ is the average density within the radius where the velocity reaches half its maximum value (Alam et al 2002). See H04 for a justification of these additional constraints, here it is sufficient to note that for the LSB galaxies $`\mathrm{\Delta }_{V/2}`$ is derived from the fit of Eq. (1) to the rotation curves, whereas $`\mathrm{\Delta }_{V/2,\mathrm{CDM}}`$ is the predicted value for a galaxy with $`V_{\mathrm{max}}=V_0`$ in the $`\mathrm{\Lambda }`$CDM model presented in Bullock et al. (2001) and Wechsler et al. (2002) (see also Fig. 2 in Alam et al. 2002 and Fig. 11 in H04). The constraint on $`\mathrm{\Delta }_{V/2}`$ ensures that only fits are considered with average densities close to those predicted in a standard $`\mathrm{\Lambda }`$CDM universe.
H04 find that for the majority ($``$ 70 per cent) of the galaxies fits can be found for which both the best-fit $`\chi _{\mathrm{red}}^2<1.5`$ \[using constraints (1)-(3)\] and the constrained $`\chi _{\mathrm{red},\mathrm{CDM}}^2<1.5`$ \[using constraints (i)-(iii)\]. This group of galaxies (Group A using the terminology in H04) is thus considered to be consistent with CDM. For the remaining galaxies no CDM-compatible fit could be found (Group B in H04), or the curves were too irregular to be well fitted by Eq. (1) (Group C in H04).
The main claim in H04 is thus that most of the LSB galaxies can still be fitted reasonably well with these extra constraints and without a large increase in $`\chi _{\mathrm{red}}^2`$. They conclude, based on the agreements of the constrained, CDM-compatible $`\gamma `$ and $`\mathrm{\Delta }_{V/2}`$ distributions with those of the simulations, that “this sample of LSB rotation curves is not manifestly inconsistent with the predictions of $`\mathrm{\Lambda }`$CDM cosmological models.” This conclusion thus contradicts most observational work on LSB rotation curves which argues the opposite.
In this paper I discuss these conclusions and investigate the method with which they are derived. I find that the method presented in H04 cannot distinguish between shallow and steep mass-density slopes, and illustrate this by showing that with this method even core-dominated, pseudo-isothermal halos would be inferred to be consistent with CDM. I introduce an additional criterion that does take the inner mass-density slope into account, and find that only a quarter of the LSB galaxies investigated are possibly consistent with CDM. However, most of these galaxies’ fit parameters are so weakly constrained that this is not a very strong conclusion. Only 3 out of the 51 galaxies investigated have tightly constrained solutions that are not significantly inconsistent with CDM. However, two of them are high surface brightness galaxies that are likely dominated by their stellar population. The net result is that in the total sample there is only one galaxy that is likely to be dark matter dominated and not significantly inconsistent with CDM. LSB galaxies that are consistent with CDM seem to be rare indeed.
In Sect. 2 I derive an analytical expression for the slope of a rotation curve, based on Eq. (1). Sect. 3 compares the slopes in both data and simulations at identical physical radii. In Sect. 4 I show that the constraints imposed by H04 are not sufficient to prove or disprove agreement with CDM. In Sect. 5 a new constraint is introduced that does take the slope into account. The results are summarised in Sect. 6.
## 2 Slopes
Equation (1) is a flexible, analytical function that can describe a wide range of rotation curves (see Courteau 1997 for more examples). One can, analogous to H04, use the fitting function to investigate the properties of the underlying data. Assuming a spherical halo dominated by dark matter we can use the inversion used in e.g. de Blok et al. (2001b) to derive the corresponding mass-density distribution:
$$4\pi G\rho (r)=2\frac{V}{r}\left(\frac{dV}{dr}\right)+\left(\frac{V}{r}\right)^2,$$
(2)
in combination with Eq. (1) yields
$$4\pi G\rho (r)=\left(\frac{V_0}{r}\right)^2\left[\frac{1+3\left(\frac{r_t}{r}\right)^\gamma }{\left(1+\left(\frac{r_t}{r}\right)^\gamma \right)^{(2+\gamma )/\gamma }}\right].$$
(3)
The logarithmic slope of the mass density distribution can now be derived from Eq. (3):
$$\alpha (r)\frac{d\mathrm{log}\rho }{d\mathrm{log}r}=\frac{2+\gamma }{1+\left(\frac{r_t}{r}\right)^\gamma }+\frac{\gamma }{1+3\left(\frac{r_t}{r}\right)^\gamma }.$$
(4)
The logarithmic, or power-law slope $`\alpha (r)`$ thus *depends on both $`\gamma `$ and $`r_t`$*. Fig. 2 shows $`\alpha (r)`$ for a number of values of $`\gamma `$ and $`r_t`$, where the values have been chosen to cover the range shown by the LSB rotation curve fits from H04. Changing $`\gamma `$ has the effect of changing the shape of the slope-radius relations; changing $`r_t`$ shifts the relations horizontally without affecting the shape for any particular $`\gamma `$. Most of the plotted slopes are remarkably *shallow* at small radii. An inner asymptotic value of $`\alpha =0`$ is predominant. The CDM slope $`\alpha 1`$ can only be reproduced with Eqs. (1) and (4) for small values of $`r_t`$ and/or extremely small values of $`\gamma `$. Similarly, at large radii the slope converges to $`\alpha =2`$. These are not naturally occurring values for $`\mathrm{\Lambda }`$CDM mass profiles, with asymptotic slopes $`\alpha 1`$ in the inner parts and $`\alpha =3`$ in the outer parts. Equation (1) is really built to mimic the same striking features of rotation curves that led to the introduction of the pseudo-isothermal (ISO) halo, namely a solid-body rise in the inner parts, and a constant (“flat”) rotation velocity in the outer parts. One needs to make very specific choices of fitting parameters to make this function resemble an NFW profile.
## 3 Comparing Observations and Simulations
The H04 simulations predict steep slopes, even at the smallest reliably resolved radii, defined by the so-called convergence radius $`r_{\mathrm{conv}}`$ (Power et al., 2003). For the H04 dwarf galaxy models $`r_{\mathrm{conv}}=0.3h^1`$ kpc or $`0.4`$ kpc for $`h=0.7`$. The slope at $`r_{\mathrm{conv}}`$ in the simulations varies between $`\alpha 1`$ and $`1.3`$.
Equation (4) can be used to derive the slopes implied by the fits to the observed LSB rotation curves. For this one needs to choose a radius at which the slopes will be evaluated. In the rest of this paper two choices will be used. Firstly, the radius of the innermost measured point of the rotation curve as given in de Blok et al. (2001b) and de Blok & Bosma (2002) will be considered<sup>1</sup><sup>1</sup>1Swaters et al. (2003) do not provide values for $`r_{\mathrm{in}}`$. Here a value of $`2^{\prime \prime }`$ is used (converted using the appropriate distance), which is the typical spacing between data points for the curves presented there.. One might argue that at $`r=r_{\mathrm{in}}`$ a different radius is probed for each rotation curve, so, secondly, a constant radius $`r=0.4`$ kpc $`r_{\mathrm{conv}}`$ kpc for all curves will be used as well. One should not expect dramatically different results from both choices though: the average value of $`r_{\mathrm{in}}`$ is $`0.45`$ kpc, very much comparable with $`r_{\mathrm{conv}}`$. It should be kept in mind that even though we evaluate the slope at these particular choices for the radius, the $`\gamma `$ and $`r_t`$ parameters (and therefore the slope) are constrained by the *entirety* of the curve.
Figure 3 shows the distribution of inferred mass density slopes derived using the LSB rotation curves and the best-fit $`\gamma `$ and $`r_t`$ parameters listed in Table 2 of H04. The slopes for both $`r=r_{\mathrm{in}}`$ and $`r=r_{\mathrm{conv}}=0.4`$ kpc are shown. It is clear that the distributions are heavily biased towards shallow slopes and very different from the values found in the simulations at these radii $`(\alpha 1)`$. The small values of $`\chi _{\mathrm{red},\mathrm{best}}^2`$ derived in H04 indicate that Eq. (1) fits the data well, and the problem is thus not with the quality of the fits. In Fig. 3 is also indicated the distribution of slopes for the galaxies that H04 claim are consistent with CDM (their Group A with $`\chi _{\mathrm{red},\mathrm{best}}^2<1.5`$ and $`\chi _{\mathrm{red},\mathrm{CDM}}^2<1.5`$). This distribution is not markedly different from that of the entire sample. From Fig. 3 one could thus conclude that the majority of slopes are inconsistent with the CDM prediction. How can this be reconciled with the H04 conclusions?
As discussed above, H04 add an extra constraint to the $`\gamma `$-parameter, and argue that the resulting agreement with the range of $`\gamma `$-values shown by the simulations implies consistency with CDM. This point will be explored in more detail below, but one can already see here that this does not solve the problem of the discrepancy between observed and simulated slopes. Consider the distribution of slopes of the subset of Group A for which $`\chi _{\mathrm{red},\mathrm{best}}^2=\chi _{\mathrm{red},\mathrm{CDM}}^2`$ (and both $`<1.5`$), i.e. galaxies for which the best fit is (apparently) already consistent with CDM. The distribution of slopes of this sub-group is also plotted in Fig. 3. The selection criteria for this sub-group favour galaxies with steeper slopes, and it should thus come as no surprise that the dominant peak at $`\alpha 0`$ has disappeared. Nevertheless it is remarkable that the distributions are still dominated by slopes $`\alpha >1`$. The distribution using $`r=r_{\mathrm{in}}`$ peaks at $`\alpha 0.25`$. The distribution using $`r=0.4`$ kpc peaks at $`\alpha 0.5`$. As we are probing the slope at $`r=0.4`$ kpc and not in the centre, we should not expect all ISO haloes to have a flat slope at this radius. The distribution of slopes found is consistent with that expected at $`r=0.4`$ kpc for an ensemble of pseudo-isothermal (ISO) halos with core-radii between $`0.5`$ and $`5`$ kpc<sup>2</sup><sup>2</sup>2Note that this range is consistent with the one derived from direct ISO rotation curve fits. For example, the 5th and 95th percentile values of the core-radius distribution for the minimum disk ISO models listed in de Blok & Bosma (2002) and McGaugh et al. (2001) are 0.6 and 5.0, respectively.. The decrease in mass-density at this radius is thus much less steep than expected for a CDM halo. The steepest slope expected for a realistic ISO halo (i.e. consistent with the smallest value $`R_C0.5`$ kpc measured in real LSB galaxies; see results in de Blok & Bosma 2002; McGaugh et al. 2001; Swaters et al. 2003) is $`\alpha 0.8`$, with the large majority of ISO haloes having slopes less steep than that. The slopes measured in the H04 simulations at $`r=0.4`$ kpc vary between $`1`$ and $`1.3`$, firmly inconsistent with the observed distributions.
The reason for the discrepancy can best be appreciated by realizing that according to Eq. (4) the slope is a function of both $`\gamma `$ and $`r_t`$. This is shown graphically in Fig. 5 which shows iso-slope contours in the $`(\gamma ,r_t)`$ plane, as evaluated at $`r=0.4`$ kpc. It is immediately obvious that $`\alpha `$ is not a unique function of $`\gamma `$, but depends equally strongly on $`r_t`$. Most combinations of $`\gamma `$ and $`r_t`$ with $`\gamma 1.5`$ and $`r_t1.5`$ yield slopes shallower than $`\alpha =0.2`$. Only a very small part of parameter space results in steep, CDM-compatible slopes (essentially only the area with $`r_t<0.6`$ and/or $`\gamma <0.2`$). Constraining $`\gamma <1`$ does not actually constrain the slope to steep, CDM-like values, but still allows values up to $`\alpha 0.2`$.
Over-plotted in Fig. 5 are the best-fit $`\gamma `$ and $`r_t`$ values of the LSB rotation curves as determined by H04. A distinction is made between the entire sample (Groups A, B and C), the “CDM-compatible” Group A, and the sub-group discussed above where the best fit is apparently already CDM-compatible. The distribution of the points is not very different from group to group, though the $`\chi _{\mathrm{red},\mathrm{best}}^2=\chi _{\mathrm{red},\mathrm{CDM}}^2`$ subgroup tends to have slightly steeper slopes, as already shown in Fig. 3. Regardless of how the total sample is divided up, the very large majority of the galaxies have slopes inconsistent with the values $`\alpha 1`$ derived in the H04 simulations. Note again that the slopes in simulations and observations are compared at *identical radii*, and that the $`\gamma `$ and $`r_t`$ values were derived using the *entirety of the rotation curves*. The extra constraints introduced by H04 are thus insufficient conditions for agreement with CDM, as constraining to $`\gamma <1`$ does not automatically imply a mass-density slope consistent with CDM. Unless one is prepared to allow for shallow slopes in CDM at these radii (making it inconsistent with its own simulations) an extra constraint on the slope is needed to ensure agreement with CDM.
## 4 Reductio ad absurdum with ISO haloes
Adding this extra constraint to the four introduced by H04 is only necessary if these prior constraints are insufficient to make a unique distinction between CDM and non-CDM models. If, for example, rotation curves known to be inconsistent with CDM are not rejected as such, the prior constraints are clearly insufficient to address CDM-related questions.
One can test the strength of these constraints by “observing” a sample of simulated ISO haloes. This will show that the H04 constraints are not sufficient as they imply that the majority of these ISO haloes, all with easily detectable shallow slopes, are consistent with CDM. ISO haloes form of course the one model known to be certainly inconsistent with CDM (if they were not, there would be no cusp/core controversy). The modeling procedure as described in de Blok et al. (2003) \[hereafter dB03\] is used. The reader is referred to that paper for an extensive description of these models. The aim there was to estimate the impact of systematic effects on the observations. For that purpose many different rotation curves were simulated and “observed” after the addition of observational uncertainties due to e.g. resolution effects and random motions. dB03 used the same datasets as the basis of their models as H04 use for their analysis, and the dB03 procedure can be directly applied here to test the fitting method and make comparisons with the H04 results.
The ISO halo parameters are chosen in such a way that the resulting curves resemble minimum disk ISO fits to the observed rotation curves. Fig. 10 from de Blok et al. (2001a) shows that mass models of LSB galaxies have asymptotic velocities $`1.8<\mathrm{log}(V_{\mathrm{}}/[\mathrm{km}\mathrm{s}^1])<2.2`$ (between $`60`$ and $`160`$ km s<sup>-1</sup>). Central dark matter core densities vary between $`2<\mathrm{log}(\rho _0/[M_{}\mathrm{pc}^3])<0`$ (between $`10^2`$ and $`1`$ $`M_{}`$ pc$`{}_{}{}^{3}]`$). The core radius $`R_C`$ follows from these two parameters, and has a range $`0.57<\mathrm{log}(R_C/[\mathrm{kpc}])<0.83`$ (between $`0.27`$ and $`6.8`$ kpc). For the models a random $`V_{\mathrm{}}`$ and $`\rho _0`$ are chosen from the respective logarithmic intervals. Note that these intervals encompass the entire observed parameter range. No account is taken of any correlations that might exist between these parameters, nor of any other uncertainties that might be present in the determination of halo parameters. This is reflected in the range of core-radii, which is somewhat larger than that found from observations (typically $`0.5R_C5`$ kpc). The choice of parameter ranges is therefore liberal, but sufficient for our purposes: this is not an attempt to model individual galaxies but a *reductio ad absurdum* proof of principle.
For each model halo an inclination, resolution and sampling interval are chosen randomly from the observed distributions described in dB03. Error-bars are assigned to each sampled point, again modeled on the observed distribution of uncertainties, with a minimum value of 4 km s<sup>-1</sup>. The curve is then corrected for inclination with a corresponding increase in the size of the error-bars. To each data point a random velocity component between $`10`$ and $`10`$ km s<sup>-1</sup> is added. This is slightly simpler than the method used in dB03, but results in similarly sized random motions. Lastly, as real rotation curves do not all extend out to the same radius, for each galaxy a random outer radius between 3 and 20 kpc is defined, out to which the galaxy is “observed”.
Figure 6 shows a small selection of the artificial ISO curves. Over-plotted is the Eq. (1) fit using the best-fit constraints (1)-(3), as given in Sec. 1. The reduced $`\chi ^2`$ and the fit parameters are also given in the figure. Figure 7 shows the distribution of the best-fit $`\gamma `$ parameters for 100 different realisations of the model. Some $`80`$ per cent of the ISO galaxies have $`\gamma 2`$. This number should be compared with H04 where it was found that $`(70\pm 5)`$ per cent of the LSB rotation curves had $`\gamma 2`$. The $`\gamma `$-distribution of LSB galaxies is at least in this respect consistent with that of ISO haloes. Fig. 8 shows a histogram of the slopes at $`r=0.4`$ kpc, both for the input ISO models and the “observed” output models, the latter derived using the best-fit $`\gamma `$ and $`r_t`$ parameters. Both the input and output slopes show a clear peak at $`\alpha =0`$ with a tail towards steeper slopes where the data are probing the edge of the dark matter core<sup>3</sup><sup>3</sup>3Note again that as we are probing the slope at $`r=0.4`$ kpc and not in the center, we do not expect to find flat slopes for all ISO haloes. The steepest ISO slopes expected are nevertheless significantly less steep than the predicted CDM values. Cf. the remarks in Sect. 3.. The output distribution has a less pronounced peak at $`\alpha =0`$, and a longer tail towards steeper slopes, but as the discussion in de Blok et al. 2001b shows, this is fully consistent with the effects of the limited resolution used to “observe” the input distribution. The range in slopes found is slightly larger than that exhibited by real LSB galaxies shown in Fig. 3. This is entirely due to the larger range in core-radii used. Fits of Eq. (1) to a set of model ISO rotation curves “observed” under the same conditions as the samples of LSB galaxies under investigation here are thus able to retain and retrieve the presence of shallow inner slopes.
In order for the H04 criteria to make a valid discrimination between CDM and non-CDM models they thus need to be able to reject the majority of the ISO models presented here. Applying these criteria to the ISO models and again fitting using Eq. (1), one finds that reducing the upper limit of $`\gamma =5`$ to $`\gamma =1`$ has only a modest effect on $`\chi _{\mathrm{red}}^2`$ and almost the entire ISO sample can still be fitted reasonably well with Eq. (1) and constraints (i)-(iii), as can be seen from the representative examples in Fig. 6. CDM constraint (iv) turns out to be the most restrictive one. It rejects about $`40`$ per cent of the simulated ISO haloes with $`\left|\mathrm{log}\mathrm{\Delta }_{V/2}\mathrm{log}\mathrm{\Delta }_{V/2,\mathrm{CDM}}\right|>0.7`$ (note that H04 show that a similar fraction of the LSB galaxies have incompatible densities as well). The bottom row of Fig. 6 shows a few of the rejected models. These do not look any different from the CDM-compatible ISO curves: condition (iv) merely removes some fraction of the ISO haloes that do not have a CDM-compatible density, without any particular selection against specific values of $`\gamma `$, $`r_t`$, $`V_0`$ or other relevant parameters. The most surprising outcome though is that of this sample of random ISO haloes almost 60 per cent match the H04 CDM criteria. A comparison of the distribution of slopes in this CDM-compatible sample of haloes with the total distribution shows that there has been no selection against shallow slopes (Fig. 9). Of the total sample of 100 observed ISO models, $`84`$ per cent has an observed slope $`\alpha >0.8`$. Once condition (iv) is applied, 57 haloes remain, of which 82 per cent have best-fit slopes $`\alpha >0.8`$. Even when the slopes are derived using the constrained “CDM compatible” fit parameters, one still finds that 70 per cent has a shallow slope $`\alpha >0.8`$. As the range in core-radii used for the models is slightly larger than found observationally, the fractions mentioned here should be considered lower limits. Regardless of how the slope is defined, shallow slopes are still profusely present (and detectable!) in the models.
A comparison with the input ISO halo parameters shows that the rejected haloes are not abnormal in any way. In Fig. 10 $`\mathrm{\Delta }_{V/2}`$ is plotted against the input ISO halo central density $`\rho _0`$. There is a clear relation, as might be expected: the average density of a halo is related to its central density. Galaxies are plotted with different symbols depending on whether they meet the $`\mathrm{\Delta }_{V/2}`$ constraint or not. It is clear that the halo models that do not meet the CDM criteria are only guilty of having average densities that are too high or too low. They are not degenerate, nor do they suffer from e.g. larger error bars. A more finely-tuned choice of halo parameters could bring the rejection rate down significantly, but as the input ranges of $`\rho _0`$ and $`V_{\mathrm{}}`$ are based on observed distributions, it is more likely that real galaxies have a wider range of densities than predicted by CDM, and condition (iv) merely selects from a much broader spectrum of galaxy properties those galaxies that happen to lie in the CDM-compatible range.
Figure 11 compares the values of $`\chi _{\mathrm{red}}^2`$ for the best-fit parameters and the CDM-compatible parameters. Haloes that were rejected by the $`\mathrm{\Delta }_{V/2}`$ criterion are indicated in the right-hand side of the plot. Following H04 one can divide the haloes in three categories: *(Group A):* rotation curves that are well fit by Eq. (1) and for which a good fit with apparently CDM-compatible parameters can also be found, i.e. $`\chi _{\mathrm{red},\mathrm{best}}^2<1.5`$ and $`\chi _{\mathrm{red},\mathrm{CDM}}^2<1.5`$. *(Group B):* rotation curves that are well fit by Eq. (1) but for which no good fit with CDM-compatible parameters can be found, usually because of a discrepant value of $`\mathrm{\Delta }_{V/2}`$; these curves have $`\chi _{\mathrm{best}}^2<1.5`$ and $`\chi _{\mathrm{red},\mathrm{CDM}}^21.5`$ or undefined (i.e., with $`\mathrm{\Delta }_{V/2}`$ out of range). *(Group C):* rotation curves that are poorly fit by Eq. (1), i.e. $`\chi _{\mathrm{red},\mathrm{best}}^21.5`$ and $`\chi _{\mathrm{red},\mathrm{CDM}}^21.5`$. The latter category only contains a few galaxies, as there are no large-scale asymmetries in the model rotation curves, and most rise smoothly and monotonically with radius. Note the numerical agreement between the observational $`\chi _{\mathrm{red}}^2`$ values presented in H04 and those of the ISO models, showing that the additional uncertainties introduced in those models are a good representation of the true observational uncertainties.
Models in Group A are, according to H04, “consistent with $`\mathrm{\Lambda }`$CDM in terms of their inferred central densities and the shape of their rotation curves.” One now has to add several caveats to this statement. Firstly, even with $`\gamma <1`$ and $`\chi _{\mathrm{red},\mathrm{CDM}}^2<1.5`$ a steep slope is not guaranteed. The constraints used to define Group A are not sufficient to filter out the shallow slopes that are still present and detectable in the ISO models, even after the addition of observational uncertainties. Secondly, as the inner slopes are well-defined and not necessarily steep, the shapes of the rotation curves are also not necessarily consistent with $`\mathrm{\Lambda }`$CDM. Constraints (i)-(iii) thus seem to have had hardly any effect in improving the agreement with CDM. The only restrictive constraint is condition (iv), the $`\mathrm{\Delta }_{V/2}`$ criterion, but even this seems to have only selected models with a certain range in average density from amongst a much larger range of models. This is illustrated in Fig. 12 where a comparison is made between average $`\mathrm{\Delta }_{V/2}`$ densities of the ISO models from Group A and B with those expected in a $`\mathrm{\Lambda }`$CDM universe. Galaxies in Groups A and B are indicated separately. In the H04 interpretation the ISO halos from Group A are consistent with CDM in terms of shape and average density, while group B is consistent in terms of shape but not density, even though Groups A and B both consist of pure ISO haloes, with easily detectable shallow slopes. The CDM constraints (i)-(iv) fail to take this into account and are thus insufficient to distinguish between CDM and non-CDM models. They cannot exclude the possibility that LSB galaxies have shallow inner mass-density slopes, even at $`r=r_{\mathrm{conv}}`$. Additional constraints are needed to make a unique identification of galaxies compatible with CDM.
## 5 Additional Constraints
Now that it is established that constraints (i)-(iv) are not sufficient to uniquely identify CDM-compatible galaxies, the next step is to identify additional conditions that can. Figures 2 and 5 show that the inner slope and therefore the shape of the rotation curve do not depend on $`\gamma `$ alone, but that there is an equally strong dependence on $`r_t`$. These Figures also explain why H04 came to the conclusion that their four constraints were sufficient: implicit in their analysis is the assumption that the shape of the rotation curve depends on $`\gamma `$ alone. While it is true that for $`r_t0.6`$ kpc a smaller $`\gamma `$ implies a steeper slope, this does not ensure a steep slope in an absolute sense. For $`\gamma <1`$, depending on the value of $`r_t`$, slopes as shallow as $`\alpha =0.2`$ are still possible. Furthermore, it is shown above that the H04 method does not take these shallow slopes into account, and thus allows models with easily detectable and prominent shallow slopes purpose-built in to pass the CDM test with flying colours.
In order to identify galaxies that are really compatible with CDM, an additional constraint is needed, in such a way that only $`(\gamma ,r_t)`$ combinations are allowed that yield a steep slope \[as well as pass conditions (i)-(iv)\]. In practice this means that $`(\gamma ,r_t)`$ parameter space needs to be searched for the minimum $`\chi ^2`$ value that still results in a steep slope<sup>4</sup><sup>4</sup>4Obviously, as Eq. (1) also depends on $`V_0`$, one needs to take that parameter into account as well. This is done by stepping through values of $`V_0`$ at fixed $`\gamma `$ and $`r_t`$, and choosing the value that yields the lowest $`\chi ^2`$. Therefore $`\chi ^2(\gamma ,r_t)=\mathrm{min}[\chi ^2(\gamma ,r_t,V_0)]`$.. In order to compare directly with the simulations the slopes are again evaluated at $`r=0.4`$ kpc. At this radius the H04 simulations show slopes $`1.3\alpha 1`$, and strictly speaking one ought to restrict the search of $`(\gamma ,r_t)`$ space to slopes $`\alpha <1`$. However, to take into account uncertainties in data and simulations, and give the CDM models as much leeway as possible, a more liberal range of $`\alpha <0.8`$ will be used. Fig. 4 shows that this is the steepest slope we can expect at this radius from an ultra-compact ISO halo with $`R_C0.5`$ kpc, and this value thus cleanly separates the ISO and CDM domains.
At this point a short digression on the significance of $`\chi _{\mathrm{red}}^2`$ is useful. In testing the goodness-of-fit of a model, the reduced $`\chi ^2`$, that is, the total $`\chi ^2`$ divided by the degrees of freedom of the data, is often used. A value of $`\chi _{\mathrm{red}}^21`$ is then taken as an indication that the model is a good description of the data. The reduced $`\chi ^2`$ is an appropriate statistic when comparing different models to a single set of data, or data with a similar number of degrees of freedom. An example would be the testing of an NFW model versus an ISO model for a single rotation curve, or a set of curves with a similar number of data points.
H04 determine for each LSB rotation curve the reduced $`\chi ^2`$ for the best-fitting $`(\gamma ,r_t)`$ model, as well as the minimum $`\chi _{\mathrm{red}}^2`$ value using their CDM constraints. They argue that for most of the galaxies the CDM constraints only cause a small increase in the *reduced* $`\chi ^2`$, and that therefore the CDM-constrained models are not much worse than the best-fit models. This is not necessarily so. Remember that the reduced $`\chi ^2`$ is the $`\chi ^2`$ divided by the degrees of freedom of the data (the number of data points minus the number of parameters in the model). It is easy to see that in order to cause the same change in *reduced* $`\chi ^2`$ the change in *absolute* $`\chi ^2`$ must be much larger for a rotation curve with a large number of data points than for a rotation curve with a small number of data points. In other words, a small change in $`\chi _{\mathrm{red}}^2`$ can be very significant (that is, unlikely to be consistent with the uncertainties in the data points) for a curve with a large number of data points, while it can be insignificant (entirely consistent with the uncertainties) for a curve with only a small number of data points.
The observed sample we are dealing with here contains a large number of rotation curves, with a wide range in the number of data points, varying from $`10`$ to over 350. The relevant parameter when judging whether a change in fit parameters yields a model that is still consistent with the data is thus *not* the reduced $`\chi ^2`$, but the *change in absolute $`\chi ^2`$*. The confidence level of any parameter combination is given by the corresponding change in $`\chi ^2`$ with respect to the minimum value. For a model with three free parameters $`(\gamma ,r_t,V_0)`$ as used here, the $`(1,2,3)\sigma `$ confidence intervals for any fit are given by $`\mathrm{\Delta }\chi ^2=(3.5,8.0,13.9)`$, respectively.
Using these criteria, a rotation curve can be defined to be consistent with CDM if the best-fit values of $`\gamma `$ and $`r_t`$ lie within $`2\sigma `$ of the closest $`(\gamma ,r_t)`$ combination with a $`\alpha <0.8`$ slope. In order for the fit to give any meaningful constraints, one should furthermore require that the area enclosed by the $`2\sigma `$ confidence level is a small fraction of the total area of parameter space investigated<sup>5</sup><sup>5</sup>5The choice for $`2\sigma `$ is not crucial. We have tested the analysis with a $`3\sigma `$ level as well. Though the number of badly constrained fits increases, the relative proportions of CDM consistent *vs* inconsistent fits does not change appreciably.. If the latter is not the case, then any conclusions regarding agreement or disagreement with CDM are only very weak at best.
Here the Group A rotation curves, defined in H04 to be consistent with their definition of CDM, are investigated again using the new constraints<sup>6</sup><sup>6</sup>6Here only the de Blok & Bosma (2002) and McGaugh et al. (2001) data are investigated. However, there is no reason to believe that the Swaters et al. (2003) data set differs from the other two in any way.. Groups B and C are already inconsistent with CDM using the more relaxed H04 constraints, and will therefore not be investigated further. For each rotation curve from Group A the minimum $`\chi ^2`$ was found using Eq. (1), and the change in $`\chi ^2`$ necessary to be consistent with a slope $`\alpha <0.8`$ was determined. The $`\sigma `$-confidence interval corresponding to this $`\mathrm{\Delta }\chi ^2`$ was also calculated, as well as the ratio $`R`$ of the area enclosed by the $`2\sigma `$ contour and the total area of $`(\gamma ,r_t)`$ parameter space searched $`(\gamma :05;r_t:05`$ kpc). As described above, the ratio $`R`$ is used to determine whether a solution is well-constrained or not. Small values of $`R`$ indicate that the $`2\sigma `$ area is only a small fraction of total parameter space, and therefore yield significant solutions. Large values of $`R`$ indicate that the fit does not constrain models in a significant way.
The results are given in Table 1, where the galaxies are sorted in order of the number of data-points in each individual rotation curve. Fig. 13 shows the $`(1,2,3)\sigma `$ confidence contours for a number of representative galaxies as described below. From studying the plots of the $`\chi ^2`$ contours for all galaxies, it was found that a value $`R=0.25`$ is a reasonable discriminator between well-constrained and badly constrained fits, but the precise value of $`R`$ is not critical (see Fig. 14 for some examples of $`R0.25`$ fits). The galaxies in Table 1 fall in four distinct categories:
1. *Well-constrained solutions, consistent with CDM:* objects with best fits within $`2\sigma `$ of a steep slope $`\alpha <0.8`$ parameter combination, and with a well-constrained $`R0.25`$ solution. An example is NGC3274, as shown in Fig. 13. This category consists of 3 out of the 36 galaxies in Group A.
2. *Well-constrained solutions, inconsistent with CDM:* objects with best fits more than $`2\sigma `$ away from a steep slope $`\alpha <0.8`$ parameter combination, and with a well-constrained $`R0.25`$ solution. An example is UGC 11819, as shown in Fig. 13. This category contains 17 out of 36 galaxies.
3. *Badly constrained solutions, consistent with CDM:* objects with best fits within $`2\sigma `$ of a steep slope $`\alpha <0.8`$ parameter combination, but with a badly constrained $`R>0.25`$ solution. An example is F563-1, as shown in Fig. 13. This category contains 12 out of 36 galaxies.
4. *Badly constrained solutions, inconsistent with CDM:* objects with best fits more than $`2\sigma `$ away from a steep slope $`\alpha <0.8`$ parameter combination, and with a badly constrained $`R>0.25`$ solution. An example is DDO 185, as shown in Fig. 13. This category contains 4 out of 36 galaxies.
It is clear that only galaxies from the first two categories, i.e. with $`R0.25`$, can help constrain halo models. A glance at Table 1 shows that most of the solutions that are consistent with CDM are actually badly constrained. This is investigated further in Fig. 15. Here $`R`$ is plotted against the number of data points in each rotation curve. A distinction is made between rotation curves that are consistent with CDM and those that are inconsistent with CDM using the constraints introduced here. Table 1 and Fig. 15 lead to several conclusions. Firstly, of the 15 galaxies that are consistent with CDM, 13 have less than 30 data points. Furthermore, of these 13 galaxies, only 1 has $`R0.25`$ (i.e. is well-constrained). Looking at the 21 galaxies that are inconsistent with CDM, only 10 have less than 30 data points. Of these 10 galaxies, all except one have $`R0.25`$. Rotation curves with more than 30 data points tend to be better constrained on average, as well as have a tendency to be inconsistent with CDM. It is also interesting to note that even amongst galaxies with only a limited number of data points, the most constrained solutions (small $`R`$) tend to be those inconsistent with CDM. The majority of solutions $`(85`$ per cent) that are consistent with CDM have only a small number of data points *and* are ill-constrained. These ill-constrained solutions therefore do not decide between any model one way or the other, and statements that these galaxies are consistent with CDM are thus not very strong. Of the well-constrained solutions the large majority is inconsistent with CDM. Furthermore, of the 13 galaxies with more than 30 data points only 2 are consistent with CDM. These two galaxies (NGC 3274 and NGC 4455) also happen to have the smallest optical scale-lengths and highest surface brightnesses from all galaxies with available photometry in the McGaugh et al. (2001), de Blok & Bosma (2002) and Swaters et al. (2003) samples: both galaxies have $`\mu _{0,B}20.0`$ mag arcsec<sup>-2</sup>. These galaxies are therefore unlikely to be dominated by dark matter in their inner parts and it is very likely that the steep slopes in these two galaxies simply reflect the stellar mass distribution (see also the discussion on NGC 3274 in Sec. 9.2.2 of de Blok & Bosma 2002).
Amongst the rotation curves with less than 30 points there is only one galaxy with $`R<0.25`$ that seems consistent with CDM. This is UGC 731, a bona-fide LSB dwarf, and is the best and only candidate in the entire observed sample for a CDM-consistent, dark-matter-dominated galaxy with a well-constrained solution. The level of agreement is not perfect though, UGC 731 is consistent with a $`\alpha =0.8`$ slope at the $`1.3\sigma `$ confidence level, with the confidence level decreasing to $`1.7\sigma `$ for an $`\alpha =1`$ slope, which is still only the upper limit of the range of slopes expected from the simulations.
Once the steep slope constraint is added to the H04 constraints there are thus only 3 galaxies left that have a well-constrained fit consistent with a steep slope and $`\mathrm{\Lambda }`$CDM. Of these 3 galaxies, 2 are of very high surface brightness and are likely to be dominated by stars in the inner parts, leaving only one possible candidate with a well-constrained CDM-compatible fit. The large majority of the well-constrained fits strongly prefer shallow slopes.
The 2$`\sigma `$ limit used here already indicates that for many galaxies it is difficult to adjust the fit parameters to get a CDM-compatible fit. Restrictive as this criterion is, it really only gives an upper limit to the true probability. While one has the freedom to adjust the parameters of an individual fit to test the likelihood of an hypothesis for one particular rotation curve, this becomes problematic when dealing with many fits. If every fit is individually adjusted to fit a particular model, does that still constitute a proper test? An analogue can perhaps be found in some of the first supernova results that indicated the existence of a cosmological constant (e.g. Riess et al. 1998). The majority of the supernovae in that early work were each individually consistent within $`2\sigma `$ with the *non-*existence of $`\mathrm{\Lambda }`$. Even though the data showed a systematic trend, each measurement could have been corrected individually to fit the no-$`\mathrm{\Lambda }`$ model. However, the joint result of these small inconsistencies has formed the basis for the discovery of the effects of $`\mathrm{\Lambda }`$ on the expansion of the universe. While it would be presumptuous to compare those results with the current analysis, it does show that it is important to pay attention to the global conclusion that is forced upon us by the data.
One thing that is clear is that galaxies that are significantly consistent with $`\mathrm{\Lambda }`$CDM in terms of their inferred central densities and the shapes of their rotation curves seem to be rare indeed.
## 6 Summary
I have investigated the claim by H04 that the majority of rotation curves of LSB galaxies are not inconsistent with CDM. I have shown that the method used is unable to distinguish between shallow and steep slopes, even when these are obviously present and easily detectable. I illustrate this by showing that with the H04 method one would infer even pseudo-isothermal, core-dominated haloes to be consistent with CDM. An extra constraint based on the mass-density slope is thus needed to make more definitive statements on any (dis)agreement with CDM. Using this extra criterion I show that only a quarter of the LSB galaxies can be said to be consistent with CDM, as opposed to the three quarters found by H04. The majority of these CDM-consistent LSB galaxies do however have fit parameters that are so ill-constrained, that they are consistent with virtually anything. Restricting the analysis to galaxies with well-constrained solutions, I find that only 3 out of the 20 well-constrained galaxies are possibly consistent with CDM; two of these are High Surface Brightness dwarf galaxies that are likely to be dominated by stars.
In summary, a comparison of the circular velocity profiles of CDM haloes with rotation curves of LSB galaxies indicates that the shapes and inferred central densities of most LSB galaxies are mostly *inconsistent* with those of simulated haloes within the limitations imposed by observational error.
It is a pleasure to thank Charles Jenkins, Stacy McGaugh, Chris Power and David Weldrake for useful discussions. I thank the anonymous referee for constructive comments.
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# Symmetries of Locally Rotationally Symmetric Models
## 1 Introduction
The purpose of this paper is to study matter collineations of locally rotationally symmetric (LRS) spacetimes. Throughout the paper $`M`$ will denote the usual smooth (connected, Hausdorff, 4-dimensional) spacetime manifold with smooth Lorentz metric $`g`$ of signature $`(+,,,)`$. Thus $`M`$ is paracompact. A comma, semi-colon and the symbol $`\mathrm{\pounds }`$ denote the usual partial, covariant and Lie derivative, respectively, the covariant derivative being with respect to the Levi-Civita connection on $`M`$ derived from $`g`$. The associated Ricci and stress-energy tensors will be denoted in component form by $`R_{ab}(R^c{}_{bcd}{}^{})`$ and $`T_{ab}`$ respectively.
In recent years, much interest has been shown in the study of matter collineation (MCs) -. A vector field along which the Lie derivative of the energy-momentum tensor vanishes is called an MC, i.e.,
$$\mathrm{\pounds }_\xi T_{ab}=0,$$
(1)
where $`\xi ^a`$ is the symmetry or collineation vector. The MC equations, in component form, can be written as
$$T_{ab,c}\xi ^c+T_{ac}\xi _{,b}^c+T_{cb}\xi _{,a}^c=0,$$
(2)
where the indices $`a,b,c`$ run from $`0`$ to $`3`$. Also, assuming the Einstein field equations, a vector $`\xi ^a`$ generates an MC if $`\mathrm{\pounds }_\xi G_{ab}=0`$. It is obvious that the symmetries of the metric tensor (isometries) are also symmetries of the Einstein tensor $`G_{ab}`$, but this is not necessarily the case for the symmetries of the Ricci tensor (Ricci collineations) which are not, in general, symmetries of the Einstein tensor. In a very recent work, M. Tsamparlis et al. have discussed Ricci and matter collineations of LRS metrics for the non-degenerate case only. Here we calculate MCs of hypersurface homogeneous spacetimes which are locally rotationally symmetric models for degenerate case and relate them with isometries.
Carot et al. and Hall et al. have noticed some important general results about the Lie algebra of MCs given in the following:
(i) The set of all MCs on $`M`$ is a vector space, but it may be infinite dimensional and may not be a Lie algebra. If $`T_{ab}`$ is degenerate, i.e., $`det(T_{ab})=0`$, then $`rank(T_{ab})<4`$ and we cannot guarantee the finite dimensionality of the MCs. If $`T_{ab}`$ is non-degenerate, i.e., $`det(T_{ab})0`$, then $`rank(T_{ab})=4`$ and the Lie algebra of MCs is finite dimensional.
(ii) If the energy-momentum tensor is of rank $`4`$ everywhere then it may be regarded as a metric on $`M`$. Then it follows by a standard result that the family of MCs is, in fact, a Lie algebra of smooth vector fields on $`M`$ of finite dimension $`10`$ (and $`9`$).
(iii) Assuming $`f`$ is a scalar function defined on $`M`$, then $`\xi =f𝐗`$ is also an MC if and only if either $`f`$ is a constant on $`M`$ or $`𝐗`$ satisfies $`T_{ab}X^a=0`$, in which case $`T_{ab}`$ is necessarily degenerate and $`𝐗`$ is an eigenvector of the energy-momentum tensor with eigenvalue $`T/2`$.
(iv) If a vector field $`\xi `$ on $`M`$ is a symmetry of all the gravitational field sources, then one could require $`\mathrm{\pounds }_\xi T_{ab}=0`$ (for the non-vacuum sources) and $`\mathrm{\pounds }_\xi C_{bcd}^a=0`$ (for the vacuum sources), where $`C_{bcd}^a`$ are Weyl curvature tensor components. This leads to a famous result given by Hall et al.
Theorem: Let $`M`$ be a spacetime manifold. Then, generically, any vector field $`\xi `$ on $`M`$ which simultaneously satisfies $`\mathrm{\pounds }_\xi T_{ab}=0`$ ($`\mathrm{\pounds }_\xi G_{ab}=0`$) and $`\mathrm{\pounds }_\xi C_{bcd}^a=0`$ is a homothetic vector field.
If $`\xi ^a`$ is a Killing vector (KV) ( or a homothetic vector), then $`\mathrm{\pounds }_\xi T_{ab}=0`$, thus every isometry is also an MC but the converse is not true, in general. Notice that collineations can be proper (non-trivial) or improper (trivial). Proper MC is defined to be an MC which is not a KV, or a homothetic vector.
The rest of the paper is organized as follows. In the next section, we write down MC equations for LRS spacetimes. In section 3, we shall solve these MC equations when the energy-momentum tensor is degenerate. In section 4, we evaluate MCs for Bianchi type V metric and finally, a summary of the results obtained will be presented.
## 2 Matter Collineation Equations
The locally rotationally symmetric spacetimes have many well known important families of exact solutions of Einstein field equations and are studied extensively -. They admit a group of motions $`G_4`$ acting multiply transitively on 3-dimensional non-null orbits spacelike ($`S_3`$) or timelike ($`T_3`$) and the isotropy group is a spatial rotation. The metrics of these models can be written in the forms
$$ds^2=ϵ[dt^2+A^2(t)dx^2]B^2(t)(dy^2+\mathrm{\Sigma }^2(y,k)dz^2),$$
(3)
$$ds^2=ϵ[dt^2+A^2(t)(dx\mathrm{\Lambda }(y,k)dz)^2]B^2(t)(dy^2+\mathrm{\Sigma }^2(y,k)dz^2),$$
(4)
$$ds^2=ϵ[dt^2+A^2(t)dx^2]B^2(t)e^{2x}(dy^2+dz^2),$$
(5)
with $`ϵ=\pm 1`$ and $`k=0,\pm 1`$, where $`\mathrm{\Sigma }(y,k)`$ is $`\mathrm{sin}y,y,\mathrm{sinh}y`$ and $`\mathrm{\Lambda }(y,k)`$ is $`\mathrm{cos}y,y^2/2,\mathrm{cosh}y`$ for $`k=+1,0,1`$ respectively. It is mentioned here that the value of $`ϵ=\pm 1`$ differentiates between the static and non-static cases as can be seen by interchanging the coordinates $`t,x`$. We restrict our attention to the non-static case ($`ϵ=1`$) as the results of the static case can be obtained consequently. The LRS metrics ($`ϵ=1`$) given by Eq.(3) turn out to be Bianchi types I (BI) or $`VII_0`$ (B$`VII_0`$) for $`k=0`$, III (BIII) for $`k=1`$ and Kantowski-Sachs (KS) for $`k=+1`$. The LRS metrics ($`ϵ=1`$) given by Eq.(4) become Bianchi types II (BII) for $`k=0`$, VIII (BVIII) or III (BIII) for $`k=1`$ and IX (BIX) for $`k=+1`$. The LRS spacetime ($`ϵ=1`$) given by Eq.(5) represents Bianchi type V(BV) or $`VII_h`$ (B$`VII_h`$) metric. A complete classification of Bianchi types I, III and Kantowski-Sachs spacetimes according to the nature of energy-momentum tensors has already been given. Here we work out MCs of BII, BVIII, BIX, and BV spacetimes only for the degenerate case.
We write down the MC equations for Bianchi types II, VIII and IX spacetimes given by Eq.(4). The non-vanishing components of Ricci and energy-momentum tensors are given in Appendix A. Using these, we can write the MC Eqs.(2) as follows
$$\dot{T}_0\xi ^0+2T_0\xi _{,0}^0=0,$$
(6)
$$\dot{T}_1\xi ^0+2T_1(\xi _{,1}^1\mathrm{\Lambda }\xi _{,1}^3)=0,$$
(7)
$$\dot{T}_2\xi ^0+2T_2\xi _{,2}^2=0,$$
(8)
$`(\mathrm{\Lambda }^2\dot{T}_1+\mathrm{\Sigma }^2\dot{T}_2)\xi ^0+2(\mathrm{\Lambda }\mathrm{\Lambda }^{}T_1+\mathrm{\Sigma }\mathrm{\Sigma }^{}T_2)\xi ^22\mathrm{\Lambda }T_1(\xi _{,3}^1`$
$`\mathrm{\Lambda }\xi _{,3}^3)+2\mathrm{\Sigma }^2T_2\xi ^3_{,3}=0,`$ (9)
$$T_0\xi _{,1}^0+T_1(\xi _{,0}^1\mathrm{\Lambda }\xi _{,0}^3)=0,$$
(10)
$$T_0\xi _{,2}^0+T_2\xi _{,0}^2=0,$$
(11)
$$T_0\xi _{,3}^0\mathrm{\Lambda }T_1(\xi _{,0}^1\mathrm{\Lambda }\xi _{,0}^3)+\mathrm{\Sigma }^2T_2\xi _{,0}^3=0,$$
(12)
$$T_1\xi _{,2}^1\mathrm{\Lambda }T_1\xi _{,2}^3+T_2\xi _{,1}^2=0,$$
(13)
$$\mathrm{\Lambda }\dot{T}_1\xi ^0+\mathrm{\Lambda }^{}T_1\xi ^2T_1(\xi _{,3}^1\mathrm{\Lambda }\xi _{,3}^3)+\mathrm{\Lambda }T_1(\xi _{,1}^1\mathrm{\Lambda }\xi _{,1}^3)\mathrm{\Sigma }^2T_2\xi _{,1}^3=0,$$
(14)
$$T_2\xi _{,3}^2\mathrm{\Lambda }T_1(\xi _{,2}^1\mathrm{\Lambda }\xi _{,2}^3)+\mathrm{\Sigma }^2T_2\xi _{,2}^3=0,$$
(15)
where dot and prime denote differentiation with respect to time coordinate ”t” and $`\mathrm{"}x\mathrm{"}`$ respectively. Notice that we have used the notation $`T_{aa}=T_a`$. We solve these equations when $`det(T_{ab})=0`$. It should be mentioned here that we shall use $`\sqrt{T_0}`$ with $`T_0>0`$ in order to fulfill the dominant energy condition.
## 3 Matter Collineations in the Degenerate Case
In order to solve MC equations (6)-(15) when at least one component of $`T_m=0,(m=0,1,2)`$, we can have the following two main cases:
(1) when only one of the $`T_m0`$,
(2) when exactly two of the $`T_m0`$,
It is mentioned here that the trivial case, where $`T_m=0`$, shows that every vector field is an MC.
Case (1): This case can further be subdivided into three cases:
(1a) $`T_00,T_j=0(j=1,2)`$,
(1b) $`T_10,T_k=0(k=0,2)`$,
(1c) $`T_20,T_l=0(l=0,1)`$.
The case (1a) is trivial and we get
$$\xi =\frac{c_0}{\sqrt{T_0}}_t+\xi ^i(x^a)_i,(i=1,2,3),$$
(16)
where $`c_0`$ is a constant.
When $`T_1=0=T_2`$, using values of these components of the energy-momentum tensor from Eqs.(A2), it follows that
$$2\frac{A^2\ddot{B}}{B}\frac{A^2\dot{B}^2}{B^2}+3\frac{A^4}{4B^4}k\frac{A^2}{B^2}=0,(k=0,+1,1),$$
(17)
$$B\ddot{B}\frac{\ddot{A}B^2}{A}\frac{\dot{A}B\dot{B}}{A}\frac{A^2}{4B^2}=0.$$
(18)
These are second order non-linear differential equations in $`A`$ and $`B`$ and can only be solved by assuming some relationship between these two functions. If we assume that $`B=cA`$, where $`c`$ is an arbitrary constant, we obtain
$`2A\ddot{A}+\dot{A}^2+{\displaystyle \frac{4kc^23}{4c^4}}=0,`$
$`2A\ddot{A}+\dot{A}^2+{\displaystyle \frac{1}{4c^4}}=0.`$ (19)
The general solution of these equations can not be found analytically. To have some particular solution, we make an assumption that $`A`$ be of the form $`A=(at+b)^n`$, where $`a,b,n`$ are arbitrary constant. We find that only the possible solution is for $`n=1`$. Thus we obtain the following solution
$$ds^2=dt^2(at+b)^2(dx\mathrm{\Lambda }(y,k)dz)^2c^2(at+b)^2(dy^2+\mathrm{\Sigma }^2(y,k)dz^2),$$
(20)
where $`4a^2c^44kc^25=0`$. It can easily be verified that these metrics represent perfect fluid dust solutions. The energy density is given by
$`\rho ={\displaystyle \frac{3a^24kc^21}{4c^4(at+b)^2}}.`$
For the case 1(b), it follows from MC Eqs.(6)-(15) that either $`T_1=constant`$ or $`\xi ^0=0`$. When $`T_1=constant`$, we get
$`\xi ^0=\xi ^0(x^a),\xi ^1=\mathrm{\Lambda }\xi ^3+C(y,z),`$
$`\xi ^2={\displaystyle \frac{1}{\mathrm{\Lambda }^{}}}C_{,3}(y,z),\xi ^3={\displaystyle \frac{1}{\mathrm{\Lambda }^{}}}C_{,2}(y,z),`$ (21)
where $`C`$ is an integration function of $`y`$ and $`z`$. It can be checked that, in this case, energy density and pressure both vanish.
In the case of 1(c), solution of the MC equations will become
$$\xi ^0=2\frac{T_2}{\dot{T}_2}c_0z\mathrm{\Sigma }^{},\xi ^1=\xi ^1(x^a),\xi ^2=c_0z\mathrm{\Sigma },\xi ^3=c_0\frac{dy}{\mathrm{\Sigma }}+c_1,$$
(22)
where $`c_0`$ and $`c_1`$ are arbitrary constants. This case also gives both pressure and energy density zero for this model. We see that the case (1) give infinite number of MCs.
Case (2): This case implies the following three possibilities:
(2a) $`T_l0,T_2=0`$,
(2b) $`T_j0,T_0=0`$,
(2c) $`T_k0,T_1=0`$.
The case 2(a) explores further two possibilities i.e. either $`T_1=constant0`$ or $`\xi ^0=0`$. For $`T_1=constant`$, solution of the MC equations yields
$$\xi =\frac{c_0}{\sqrt{T_0}}_t+\xi ^1(z)_x+\frac{1}{\mathrm{\Lambda }^{}}\xi _{,3}^1_y.$$
(23)
When $`T_2=0`$, using the same procedure as in the case 1(a), we find the same solution with the condition given as $`4a^2c^4+1=0`$.
In the case 2(c), from MC Eqs.(6)-(15), in addition to the improper MCs $`\xi _{(1)},\xi _{(2)},\xi _{(3)}`$ given in Appendix A, we obtain the following proper MCs
$`\xi _{(4)}=\xi ^1(x^a)_x,`$
$`\xi _{(5)}={\displaystyle \frac{\mathrm{\Sigma }\mathrm{sin}z}{\sqrt{T_0}}}[{\displaystyle \frac{\dot{T_2}}{2kT_2}}𝐗_1+_t],`$
$`\xi _{(6)}={\displaystyle \frac{\mathrm{\Sigma }\mathrm{cos}z}{\sqrt{T_0}}}[{\displaystyle \frac{\dot{T_2}}{2kT_2}}𝐗_2+_t],`$
$`\xi _{(7)}={\displaystyle \frac{\mathrm{\Sigma }}{\sqrt{T_0}}}[{\displaystyle \frac{\dot{T_2}}{2T_2}}_y+{\displaystyle \frac{\mathrm{\Sigma }^{}}{\mathrm{\Sigma }}}_t],`$ (24)
in which the constraint equation is
$$\frac{T_2}{\sqrt{T_0}}[\frac{\dot{T_2}}{2T_2\sqrt{T_0}}\dot{]}=k.$$
(25)
The values of $`𝐗_1`$ and $`𝐗_2`$ are given as
$`𝐗_1={\displaystyle \frac{\mathrm{\Sigma }^{}}{\mathrm{\Sigma }}}_y+{\displaystyle \frac{\mathrm{cot}z}{\mathrm{\Sigma }^2}}_z,`$
$`𝐗_2={\displaystyle \frac{\mathrm{\Sigma }^{}}{\mathrm{\Sigma }}}_y+{\displaystyle \frac{\mathrm{tan}z}{\mathrm{\Sigma }^2}}_z,`$ (26)
Again we see that all the possibilities of the case (2) give infinite-dimensional MCs.
## 4 Matter Collineations of Bianchi Type V Metric
The LRS spacetime ($`ϵ=1`$) given by Eq.(5) represents Bianchi type V metric. MC equations for this metric will become
$$T_{0,0}\xi ^0+2T_0\xi _{,0}^0=0,$$
(27)
$$T_{1,0}\xi ^0+2T_1\xi _{,1}^1=0,$$
(28)
$$\xi _{,2}^2\xi _{,3}^3=0,(T_20),$$
(29)
$$T_0\xi _{,1}^0+T_1\xi _{,0}^1=0,$$
(30)
$$T_0\xi _{,2}^0+T_2\xi _{,0}^2=0,$$
(31)
$$T_0\xi _{,3}^0+T_2\xi _{,0}^3=0,$$
(32)
$$T_1\xi _{,2}^1+T_2\xi _{,1}^2=0,$$
(33)
$$T_1\xi _{,3}^1T_2\xi _{,1}^3=0,$$
(34)
$$\xi _{,3}^2+\xi _{,2}^3=0,(T_20).$$
(35)
Again we shall restrict ourselves only for the degenerate case. The non-vanishing components of Ricci and energy-momentum tensor are given in Appendix B. There arises two possibilities for this case according as (1) when one of the components of energy-momentum tensor is non-zero and (2) when two of the components are non-zero. As we have done previously, each of these two cases can further be divided into three subcases.
The case 1(a) yields the same solution as for the case 1(a) of the previous section. Using the Einstein field equations for the perfect fluid matter, we find that the model is pressureless (i.e. $`p=0`$) and energy density is given as
$`\rho =2{\displaystyle \frac{\dot{A}\dot{B}}{AB}}+{\displaystyle \frac{\dot{B}^2}{B^2}}{\displaystyle \frac{1}{A^2}}.`$
For the case (1b), it follows from MC equations that
$$\xi =\frac{c_0}{\sqrt{T_1}}_x+\xi ^n(x^a)_n,(n=0,2,3),$$
(36)
where $`c_0`$ is a constant.
In the case (1c), we obtain the following solution
$$\xi =\xi ^l(x^a)_l+c_0_y+c_1_z.$$
(37)
It is obvious that all the possibilities of the case (1) give infinite dimensional MCs.
For the case 2(a), when we solve MC equations, we obtain the following constraint
$$\frac{T_1}{\sqrt{T_0}}(\frac{\dot{T}_1}{2T_1\sqrt{T_0}}\dot{)}=\alpha ,$$
(38)
where $`\alpha `$ is an arbitrary constant which can be positive, zero or negative.
When $`\alpha >0`$, we have the following solution
$`\xi =c_0_x+c_1{\displaystyle \frac{1}{\sqrt{T_0}}}(\mathrm{cos}\sqrt{\alpha }x_t{\displaystyle \frac{\dot{T}_1}{2T_1\sqrt{\alpha }}}\mathrm{sin}\sqrt{\alpha }x_x)`$
$`+c_2{\displaystyle \frac{1}{\sqrt{T_0}}}(\mathrm{sin}\sqrt{\alpha }x_t+{\displaystyle \frac{\dot{T}_1}{2T_1\sqrt{\alpha }}}\mathrm{cos}\sqrt{\alpha }x_x)`$
$`+\xi ^2(x^a)_y+\xi ^3(x^a)_z.`$ (39)
The case $`\alpha =0`$ yields the solution
$`\xi =c_0_x+c_1{\displaystyle \frac{1}{\sqrt{T_0}}}(x_t({\displaystyle \frac{\dot{T}_1}{4T_1}}x^2+\sqrt{T_0}{\displaystyle \frac{\sqrt{T_0}}{T_1}𝑑t})_x)`$
$`+c_2{\displaystyle \frac{1}{\sqrt{T_0}}}(_t{\displaystyle \frac{\dot{T}_1}{2T_1}}x_x)+\xi ^2(x^a)_y+\xi ^3(x^a)_z.`$ (40)
For $`\alpha <0`$, we obtain
$`\xi =c_0_x+c_1{\displaystyle \frac{1}{\sqrt{T_0}}}(\mathrm{cosh}\sqrt{\alpha }x_t{\displaystyle \frac{1}{2\sqrt{\alpha }T_1}}(\dot{T}_1+4\alpha T_0)\mathrm{sinh}\sqrt{\alpha }x_x)`$
$`+c_2{\displaystyle \frac{1}{\sqrt{T_0}}}(\mathrm{sinh}\sqrt{\alpha }x_t{\displaystyle \frac{1}{2\sqrt{\alpha }T_1}}(\dot{T}_1+4\alpha T_0)\mathrm{cosh}\sqrt{\alpha }x_x)`$
$`+\xi ^2(x^a)_y+\xi ^3(x^a)_z.`$ (41)
Again we see that we obtain infinite dimensional MCs.
In the case of 2(b), solving MC equations simultaneously, we obtain
$$\xi =2\frac{T_1}{\dot{T_1}}A^{}(x)_t+A(x)_x+c_0(_y+_z)+c_1\frac{1}{2}(z_yy_z).$$
(42)
When we solve MC equations using the constraints of the case 2(c), it follows that
$$\xi =\frac{c_0}{\sqrt{T_0}}_t+\xi ^1(x^a)_x+c_1(_y+_z)+c_2\frac{1}{2}(z_yy_z).$$
(43)
We see that $`\xi ^1`$ is an arbitrary function depending on all four variables, thus we have infinite dimensional MCs.
## 5 Conclusion
This paper has been devoted to the evaluation of MCs for the LRS models when the energy-momentum tensor is degenerate. We have concentrated only for Bianchi types II, VIII, IX and IX spacetimes given by Eqs.(4) and (5) respectively as the classification of metrics given by Eq.(3) has already been completed . It is found that when at least one component of $`T_{ab}`$ is non-zero (case (1) of sec. 3), all the possibilities yield infinite dimensional MCs. If one component of $`T_{ab}`$ is zero (case (2) of sec. 3), we again obtain infinite-dimensional MCs. However, in this case, we have proper MCs, in addition, to the improper MCs. For this case, the solution of the equation $`T_2=0`$ turns out to be the perfect fluid dust solution. We also observe that all cases of the Bianchi type V yield infinite dimensional MCs.
We note that the case in which $`T_00(T_0>0)`$ is the only surviving component of $`T_{ab}`$ can always be interpreted as a dust fluid. In the case when $`T_0=0`$, we do not have dominant energy condition instead we have energy density zero.
As we have considered degenerate case, it is natural to expect infinite dimensional MCs as given by Hall et al. . However, it may be finite dimensional for some case as given in when we classify MCs for the LRS models given by Eq.(3). The case 2(b) (section 3) is left open and it needs to be investigated. It is expected that this case would provide finite dimensional MCs. Also, we have obtained constraints on the energy-momentum tensor. It might be interesting to look for more solutions of the constraint equations or at least examples should be constructed to satisfy these constraints.
## Appendix A
The surviving components of the Ricci tensor are
$`R_{00}`$ $`=`$ $`{\displaystyle \frac{1}{AB}}(2A\ddot{B}+\ddot{A}B),`$
$`R_{11}`$ $`=`$ $`A\ddot{A}+2{\displaystyle \frac{A\dot{A}\dot{B}}{B}}+{\displaystyle \frac{A^4}{2B^4}},`$
$`R_{22}`$ $`=`$ $`B\ddot{B}+\dot{B}^2{\displaystyle \frac{A^2}{2B^2}}+{\displaystyle \frac{\dot{A}\dot{B}B}{A}}+k,`$
$`R_{33}`$ $`=`$ $`\mathrm{\Lambda }^2R_{11}+\mathrm{\Sigma }^2R_{22},`$
$`R_{13}`$ $`=`$ $`\mathrm{\Lambda }R_{11}.`$ (A1)
where dot represents derivative w.r.t. time coordinate $`t`$. The non-vanishing components of energy-momentum tensor $`T_{ab}`$ are
$`T_{00}`$ $`=`$ $`2{\displaystyle \frac{\dot{A}\dot{B}}{AB}}+{\displaystyle \frac{\dot{B}^2}{B^2}}{\displaystyle \frac{A^2}{4B^4}}+{\displaystyle \frac{k}{B^2}},`$
$`T_{11}`$ $`=`$ $`2{\displaystyle \frac{A^2\ddot{B}}{B}}{\displaystyle \frac{A^2\dot{B}^2}{B^2}}+3{\displaystyle \frac{A^4}{4B^4}}k{\displaystyle \frac{A^2}{B^2}},`$
$`T_{22}`$ $`=`$ $`B\ddot{B}{\displaystyle \frac{\ddot{A}B^2}{A}}{\displaystyle \frac{\dot{A}B\dot{B}}{A}}{\displaystyle \frac{A^2}{4B^2}},`$
$`T_{33}`$ $`=`$ $`\mathrm{\Lambda }^2T_{11}+\mathrm{\Sigma }^2T_{22},`$
$`T_{13}`$ $`=`$ $`\mathrm{\Lambda }T_{11}.`$ (A2)
The independent KVs associated with the LRS spacetimes are given by
$`\xi _{(1)}`$ $`=`$ $`\mathrm{sin}\varphi _\theta +\mathrm{cot}\theta \mathrm{cos}\varphi _\varphi ,`$
$`\xi _{(2)}`$ $`=`$ $`\mathrm{cos}\varphi _\theta \mathrm{cot}\theta \mathrm{sin}\varphi _\varphi ,`$
$`\xi _{(3)}`$ $`=`$ $`_\varphi .`$ (A3)
## Appendix B
The surviving components of the Ricci tensor are
$`R_0`$ $`=`$ $`{\displaystyle \frac{\ddot{A}}{A}}2{\displaystyle \frac{\ddot{B}}{B}},`$
$`R_1`$ $`=`$ $`A\ddot{A}+2{\displaystyle \frac{A\dot{A}\dot{B}}{B}},`$
$`R_2`$ $`=`$ $`(B\ddot{B}+\dot{B}^2{\displaystyle \frac{B^2}{A^2}}+{\displaystyle \frac{\dot{A}\dot{B}B}{A}}=R_3.`$ (B1)
The non-vanishing components of energy-momentum tensor $`T_{ab}`$ are
$`T_0`$ $`=`$ $`2{\displaystyle \frac{\dot{A}\dot{B}}{AB}}+{\displaystyle \frac{\dot{B}^2}{B^2}}{\displaystyle \frac{1}{A^2}},`$
$`T_1`$ $`=`$ $`2{\displaystyle \frac{A^2\ddot{B}}{B}}{\displaystyle \frac{A^2\dot{B}^2}{B^2}}+1,`$
$`T_2`$ $`=`$ $`(B\ddot{B}+{\displaystyle \frac{\ddot{A}B^2}{A}}+{\displaystyle \frac{\dot{A}\dot{B}B}{A}})e^{2x}=T_3.`$ (B2)
Acknowledgment
The author would like to thank the referee for useful and accurate remarks and suggestions.
References
Carot, J., da Costa, J. and Vaz, E.G.L.R.: J. Math. Phys. 35(1994)4832.
Hall, G.S., Roy, I. and Vaz, L.R.: Gen. Rel and Grav. 28(1996)299.
Carot, J. and da Costa, J.: Procs. of the 6th Canadian Conf. on General Relativity and Relativistic Astrophysics, Fields Inst. Commun. 15, Amer. Math. Soc. WC Providence, RI(1997)179.
Yavuz, İ., and Camcı, U.: Gen. Rel. Grav.28(1996)691;
Camcı, U., Yavuz, İ., Baysal, H., Tarhan, İ., and Yılmaz, İ.: Int. J. Mod. Phys. D10(2001)751.
Camcı, U. and Barnes, A.: Class. Quant. Grav. 19(2002)393.
Sharif, M.: Nuovo Cimento B116(2001)673; Astrophys. Space Sci. 278(2001)447; J. Math. Phys. (2003);
Sharif, M. and Sehar Aziz: Gen Rel. and Grav. 35(2003).
Camcı, U. and Sharif, M.: Gen Rel. and Grav. 35(2003)97.
Camcı, U. and Sharif, M.: Class. Quant. Grav. 20(2003)2169.
Tsamparlis, M. and Apostolopoulos, Pantelis S.: Gen. Rel. and Grav. 36(2004)47.
Ellis, G.F.R.: J. Math. Phys. 8(1967)1171.
Stewart, J.M. and Ellis, G.F.R.: J. Math. Phys. 9(1968)1072.
Ellis, G.F.R. and MacCallum, M.A.H.: Commun. Math. Phys. 12(1969)108;
MacCallum, M.A.H.: Commun. Math. Phys. 20(1971)57;
Collins, C.B.: Commun. Math. Phys. 23(1971)137.
van Elst, H. and Ellis, G.F.R.: Class. Quantum Grav. 13(1996)1099.
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# Untitled Document
IFT–05/13
The simplest 3+2 model
with two light sterile neutrinos
Wojciech Królikowski
Institute of Theoretical Physics, Warsaw University
Hoża 69, PL–00–681 Warszawa, Poland
Abstract
The simplest 3+2 neutrino model is described, where two light sterile neutrinos mix very weakly with three active neutrinos and mutually do not mix at all, while the mass-squared scale of the possible LSND effect is provided by $`m_5^2m_4^2`$ ($`\nu _4`$ and $`\nu _5`$ being two additional mass neutrinos connected with the existence of two sterile neutrinos $`\nu _s`$ and $`\nu _s^{}`$). This 3+2 model is not better for explaining the LSND effect than the simplest 3+1 neutrino model, where one light sterile neutrino mixes very weakly with three active neutrinos, while the mass-squared scale of the possible LSND effect is given by $`m_4^2m_1^2`$ ($`\nu _4`$ denoting an additional mass neutrino existing due to a sterile neutrino $`\nu _s`$). However, a small LSND effect with amplitude of the order O($`10^3`$) is not excluded by the present (pre-MiniBooNE) data.
PACS numbers: 12.15.Ff , 14.60.Pq , 12.15.Hh .
June 2005
As is well known, the neutrino mixing matrix $`U^{(3)}=\left(U_{\alpha i}^{(3)}\right)(\alpha =e,\mu ,\tau \mathrm{and}i=1,2,3)`$ appearing in the unitary transformation
$$\nu _\alpha =\underset{i}{}U_{\alpha i}^{(3)}\nu _i$$
(1)
between the flavor neutrinos $`\nu _e,\nu _\mu ,\nu _\tau `$ and mass neutrinos $`\nu _1,\nu _2,\nu _3`$ is experimentally consistent with the global bilarge form
$$U^{(3)}=\left(\begin{array}{ccc}c_{12}& s_{12}& 0\\ \frac{1}{\sqrt{2}}s_{12}& \frac{1}{\sqrt{2}}c_{12}& \frac{1}{\sqrt{2}}\\ \frac{1}{\sqrt{2}}s_{12}& \frac{1}{\sqrt{2}}c_{12}& \frac{1}{\sqrt{2}}\end{array}\right),$$
(2)
where $`s_{12}^20.30`$ , while $`U_{e3}^{(3)}=s_{13}\mathrm{exp}(i\delta )`$ is negligible according to the negative result of Chooz experiment (the upper limit is $`s_{13}^2<0.03`$). When neglecting $`s_{13}`$, the neutrino oscillation probabilities (in the vacuum) are
$$P(\nu _\alpha \nu _\beta )=\delta _{\beta \alpha }4\underset{j>i}{}U_{\beta j}^{(3)}U_{\alpha j}^{(3)}U_{\beta i}^{(3)}U_{\alpha i}^{(3)}\mathrm{sin}^2x_{ji},$$
(3)
where
$$x_{ji}1.27\frac{\mathrm{\Delta }m_{ji}^2L}{E},\mathrm{\Delta }m_{ji}^2m_j^2m_i^2$$
(4)
($`\mathrm{\Delta }m_{ji}^2`$, $`L`$ and $`E`$ are measured in eV<sup>2</sup>, km and GeV, respectively), giving experimentally $`\mathrm{\Delta }m_{21}^28.0\times 10^5\mathrm{eV}^2`$ as well as $`|\mathrm{\Delta }m_{32}^2|2.2\times 10^3\mathrm{eV}^2`$ for solar $`\nu _e`$’s and KamLAND reactor $`\overline{\nu }_e`$’s (with the solar MSW effect included) as well as atmospheric and K2K accelerator $`\nu _\mu `$’s, respectively.
The formulae (2) and (3) are consistent with the zero LSND effect. The nonzero LSND effect would require the existence of a third neutrino mass-squared splitting, absent in the case of three active neutrinos only (unless the CPT invariance of neutrino oscillations is seriously violated, what does not seem to be realistic). The LSND effect will be tested soon in the ongoing MiniBooNE experiment . If this test confirms the LSND effect, we will need at least one light sterile neutrino in addition to three active neutrinos in order to introduce one extra mass-squared splitting (but – at the same time – not to change significantly the fit to solar, reactor, atmospheric and accelerator neutrino experiments).
While the 3+1 neutrino models with one light sterile neutrino are considered to be disfavored by present data , the 3+2 neutrino schemes with two light sterile neutrinos may provide a better description of current neutrino oscillations including the LSND effect . However, the special 3+2 model of Ref. , where among three active neutrinos $`\nu _e,\nu _\mu ,\nu _\tau `$ and two sterile neutrinos $`\nu _s,\nu _s^{}`$ there are two maximally mixing pairs $`\nu _\mu ,\nu _\tau `$ and $`\nu _s,\nu _s^{}`$, does not meet these expectations. Also, the simplest 3+2 neutrino model considered in the present note is not better for explaining the LSND effect than the simplest 3+1 neutrino model .
For the experimental existence of two light sterile neutrinos (beside three generations of SM-active leptons and quarks) we argued some time ago in Refs. , where a Pauli principle working intrinsically within a generalized Dirac equation was introduced. In our simplest 3+2 neutrino model considered here, two light sterile neutrinos $`\nu _s,\nu _s^{}`$ mix very weakly with three active neutrinos $`\nu _e,\nu _\mu ,\nu _\tau `$ and mutually do not mix at all.
The simplest 3+2 neutrino model is defined by the $`5\times 5`$ mixing matrix $`U^{(5)}=\left(U_{\alpha i}^{(5)}\right)`$ ($`\alpha =e,\mu ,\tau ,s,s^{}`$ and $`i=1,2,3,4,5`$) of the form
$`U^{(5)}`$ $`=`$ $`U^{(5)}(12)U^{(5)}(14,25)`$ (10)
$`=`$ $`\left(\begin{array}{ccccc}c_{12}c_{14}& s_{12}c_{25}& 0& c_{12}s_{14}& s_{12}s_{25}\\ \frac{1}{\sqrt{2}}s_{12}c_{14}& \frac{1}{\sqrt{2}}c_{12}c_{25}& \frac{1}{\sqrt{2}}& \frac{1}{\sqrt{2}}s_{12}s_{14}& \frac{1}{\sqrt{2}}c_{12}s_{25}\\ \frac{1}{\sqrt{2}}s_{12}c_{14}& \frac{1}{\sqrt{2}}c_{12}c_{25}& \frac{1}{\sqrt{2}}& \frac{1}{\sqrt{2}}s_{12}s_{14}& \frac{1}{\sqrt{2}}c_{12}s_{25}\\ s_{14}& 0& 0& c_{14}& 0\\ 0& s_{25}& 0& 0& c_{25}\end{array}\right),`$
where
$`U^{(5)}(12)`$ $`=`$ $`\left(\begin{array}{ccccc}c_{12}& s_{12}& 0& 0& 0\\ \frac{1}{\sqrt{2}}s_{12}& \frac{1}{\sqrt{2}}c_{12}& \frac{1}{\sqrt{2}}& 0& 0\\ \frac{1}{\sqrt{2}}s_{12}& \frac{1}{\sqrt{2}}c_{12}& \frac{1}{\sqrt{2}}& 0& 0\\ 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 1\end{array}\right),`$ (16)
$`U^{(5)}(14,25)`$ $`=`$ $`\left(\begin{array}{ccccc}c_{14}& 0& 0& s_{14}& 0\\ 0& c_{25}& 0& 0& s_{25}\\ 0& 0& 1& 0& 0\\ s_{14}& 0& 0& c_{14}& 0\\ 0& s_{25}& 0& 0& c_{25}\end{array}\right).`$ (22)
In this model, the extended unitary transformation (1) holds between the flavor neutrinos $`\nu _e,\nu _\mu ,\nu _\tau ,\nu _s,\nu _s^{}`$ and mass neutrinos $`\nu _1,\nu _2,\nu _3,\nu _4,\nu _5`$, where $`U_{\alpha i}^{(3)}`$ are replaced by $`U_{\alpha i}^{(5)}`$. Here, $`c_{14}^2s_{14}^2`$ and $`c_{25}^2s_{25}^2`$ guarantee the very weak mixing of two sterile neutrinos $`\nu _s,\nu _s^{}`$ with three active neutrinos $`\nu _e,\nu _\mu ,\nu _\tau `$. Besides, there is no mixing between $`\nu _s`$ and $`\nu _s^{}`$.
Then, making use of the extended neutrino oscillation formulae (3), where $`U_{\alpha i}^{(3)}`$ are replaced by $`U_{\alpha i}^{(5)}`$, we obtain (in the vacuum)
$$P(\nu _e\nu _e)14c_{12}^2s_{12}^2\left(1s_{14}^2s_{25}^2\right)\mathrm{sin}^2x_{21}4c_{12}^2s_{14}^2\mathrm{sin}^2x_{41}4s_{12}^2s_{25}^2\mathrm{sin}^2x_{51}$$
(23)
when $`x_{41}x_{42}`$, $`x_{51}x_{52}`$, and
$`P(\nu _\mu \nu _\mu )`$ $``$ $`1c_{12}^2s_{12}^2\left(1s_{14}^2s_{25}^2\right)\mathrm{sin}^2x_{21}\left(1s_{12}^2s_{14}^2c_{12}^2s_{25}^2\right)\mathrm{sin}^2x_{32}`$ (24)
$`2s_{12}^2s_{14}^2\mathrm{sin}^2x_{41}2c_{12}^2s_{25}^2\mathrm{sin}^2x_{51}`$
when $`x_{31}x_{32}`$, $`x_{41}x_{42}`$, $`x_{51}x_{52}`$, as well as
$`P(\overline{\nu }_\mu \overline{\nu }_e)`$ $``$ $`2c_{12}^2s_{12}^2\left(1s_{14}^2s_{25}^2\right)\mathrm{sin}^2x_{21}+2c_{12}^2s_{12}^2s_{14}^2s_{25}^2\mathrm{sin}^2x_{54}`$ (25)
$`+2c_{12}^2s_{12}^2\left(s_{14}^2s_{25}^2\right)\left(s_{14}^2\mathrm{sin}^2x_{41}s_{25}^2\mathrm{sin}^2x_{51}\right)`$
when $`x_{41}x_{42}`$, $`x_{51}x_{52}`$. In Eqs. (7) and (8), the terms $`O(s_{14}^4)`$, $`O(s_{25}^4)`$ and $`O(s_{14}^2s_{25}^2)`$ are neglected.
\>From Eqs. (7) and (8) as well as (9) we get
$`P(\nu _e\nu _e)_{\mathrm{sol}}`$ $``$ $`14c_{12}^2s_{12}^2(1s_{14}^2s_{25}^2)\mathrm{sin}^2(x_{21})_{\mathrm{sol}}2(c_{12}^2s_{14}^2+s_{12}^2s_{25}^2)`$ (26)
$`=`$ $`(1s_{14}^2s_{25}^2)[14c_{12}^2s_{12}^2\mathrm{sin}^2(x_{21})_{\mathrm{sol}}](c_{12}^2s_{12}^2)(s_{14}^2s_{25}^2)`$
when $`x_{21}x_{41}<x_{51}`$ with $`(x_{21})_{\mathrm{sol}}=O(\pi /2)`$ or
$$P(\overline{\nu }_e\overline{\nu }_e)_{\mathrm{Chooz}}12\left(c_{12}^2s_{14}^2+s_{12}^2s_{25}^2\right)1(\mathrm{experimentally})$$
(27)
when $`x_{21}x_{41}<x_{51}`$ with $`(x_{31})_{\mathrm{Chooz}}(x_{31})_{\mathrm{atm}}=O(\pi /2)`$, and
$$P(\nu _\mu \nu _\mu )_{\mathrm{atm}}(1s_{12}^2s_{14}^2c_{12}^2s_{25}^2)[1\mathrm{sin}^2(x_{32})_{\mathrm{atm}}]$$
(28)
when $`x_{21}x_{31}x_{41}<x_{51}`$ with $`(x_{31})_{\mathrm{atm}}=O(\pi /2)`$, as well as
$$P(\overline{\nu }_\mu \overline{\nu }_e)_{\mathrm{LSND}}2c_{12}^2s_{12}^2[s_{14}^2s_{25}^2\mathrm{sin}^2(x_{54})_{\mathrm{LSND}}+\frac{1}{2}(s_{14}^2s_{25}^2)^2]$$
(29)
when $`x_{21}x_{54}x_{41}<x_{51}`$ with $`(x_{54})_{\mathrm{LSND}}=O(\pi /2)`$. In the symmetric case of $`s_{14}^2s_{25}^2`$, Eq. (13) reads
$$P(\overline{\nu }_\mu \overline{\nu }_e)_{\mathrm{LSND}}2c_{12}^2s_{12}^2s_{14}^4\mathrm{sin}^2(x_{54})_{\mathrm{LSND}}.$$
(30)
If the nonzero LSND effect exists in the order $`P(\overline{\nu }_\mu \overline{\nu }_e)_{\mathrm{LSND}}(10^2\mathrm{to}\mathrm{\hspace{0.33em}10}^3)`$ $`\times \mathrm{sin}^2(x_{54})_{\mathrm{LSND}}`$, the formula (14) valid in the case of $`s_{14}^2s_{25}^2`$ gives
$$s_{14}^2\left(\frac{10^2\mathrm{to}\mathrm{\hspace{0.33em}10}^3}{2c_{12}^2s_{12}^2}\right)^{1/2}0.15\mathrm{to}\mathrm{\hspace{0.33em}0.049},$$
(31)
where $`2c_{12}^2s_{12}^20.42`$ ($`s_{12}^20.30`$). Then, in the case of $`s_{12}^2s_{25}^2`$
$$P(\nu _e\nu _e)_{\mathrm{sol}}[1(0.31\mathrm{to}\mathrm{\hspace{0.33em}0.098})][10.84\mathrm{sin}^2(x_{21})_{\mathrm{sol}}]$$
(32)
or
$$P(\overline{\nu }_e\overline{\nu }_e)_{\mathrm{Chooz}}1(0.31\mathrm{to}\mathrm{\hspace{0.33em}0.098})1(\mathrm{experimentally}),$$
(33)
and
$$P(\nu _\mu \nu _\mu )_{\mathrm{atm}}[1(0.15\mathrm{to}\mathrm{\hspace{0.33em}0.049})][1\mathrm{sin}^2(x_{32})_{\mathrm{atm}}].$$
(34)
Here, $`\mathrm{\Delta }m_{54}^2\mathrm{\Delta }m_{\mathrm{LSND}}^21\mathrm{eV}^2`$ (say).
We can see that, in particular, the nonzero LSND effect of the order $`P(\overline{\nu }_\mu \overline{\nu }_e)10^2\mathrm{sin}^2(x_{54})_{\mathrm{LSND}}`$ would imply an experimentally visible Chooz effect $`P(\overline{\nu }_e\overline{\nu }_e)10.31<1`$, which is not observed. If its order was $`P(\overline{\nu }_\mu \overline{\nu }_e)10^3\mathrm{sin}^2(x_{54})_{\mathrm{LSND}}`$ the corresponding Chooz effect, $`P(\overline{\nu }_e\overline{\nu }_e)10.0981`$, would be nearly at the (Chooz) experimental edge (here, still $`s_{13}=0`$).
In order to pass to the simplest 3+1 neutrino model , we can put $`s_{25}0`$ and $`(x_{41})_{\mathrm{LSND}}=O(\pi /2)`$, what gives
$$P(\nu _e\nu _e)_{\mathrm{sol}}(1s_{14}^2)[14c_{12}^2s_{12}^2\mathrm{sin}^2(x_{21})_{\mathrm{sol}}](c_{12}^2s_{12}^2)s_{14}^2$$
(35)
or
$$P(\overline{\nu }_e\overline{\nu }_e)_{\mathrm{Chooz}}12c_{12}^2s_{14}^21(\mathrm{experimentally}),$$
(36)
and
$$P(\nu _\mu \nu _\mu )_{\mathrm{atm}}(1s_{12}^2s_{14}^2)[1\mathrm{sin}^2(x_{32})_{\mathrm{atm}}],$$
(37)
as well as
$$P(\overline{\nu }_\mu \overline{\nu }_e)_{\mathrm{LSND}}2c_{12}^2s_{12}^2s_{14}^4\mathrm{sin}^2(x_{41})_{\mathrm{LSND}}$$
(38)
(in Eqs (19), (20) and (21) the terms $`O(s_{14}^4)`$ are neglected). Then, as before, $`s_{14}^20.15`$ to 0.049 for the nonzero LSND effect with the amplitude of the order $`O(10^2\mathrm{to}\mathrm{\hspace{0.33em}10}^3)`$. Here, $`\mathrm{\Delta }m_{41}^2\mathrm{\Delta }m_{\mathrm{LSND}}^21\mathrm{eV}^2`$ (say). Thus, the LSND effect is here the same as in our simplest 3+2 neutrino model with $`s_{14}^2s_{25}^2`$ and $`(x_{54})_{\mathrm{LSND}}=O(\pi /2)`$. But there, $`\mathrm{\Delta }m_{54}^2\mathrm{\Delta }m_{\mathrm{LSND}}^21\mathrm{eV}^2`$ (say). In both cases, the present (pre-MiniBooNE) data do not exclude a small LSND effect with the amplitude of the order $`O(10^3)`$.
References
For a review cf. S. Pascoli, S.T. Petcov and T. Schwetz, hep--ph/0505226.
M. Apollonio et al. (Chooz Collaboration), Eur. Phys. J. C 27, 331 (2003).
C. Athanassopoulos et al. (LSND Collaboration), Phys. Rev. Lett. 77, 3082 (1996); Phys. Rev. C 58, 2489 (1998); A. Aguilar et al., Phys. Rev. D 64, 112007 (2001).
A.O. Bazarko et al. (MiniBooNE Collaboration), hep--ex/9906003.
Cf. e.g. M. Maltoni, T. Schwetz, M.A. Tortola and J.W. Valle, Nucl. Phys. B 643, 321 (2002).
M. Sorel, J. Conrad and M. Shaevitz, hep--ph/0305255.
W. Królikowski, Acta Phys. Pol. B 35, 1675 (2004) \[hep--ph/0402183\].
Ref. , Appendix A.
W. Królikowski, in Proc. of the 12th Max Born Symposium, Wrocław, Poland, 1998, eds. A. Borowiec et al., Springer, 2000 \[hep-ph/9808307\]; Acta Phys. Pol. B 30, 227 (1999) \[hep-ph/9808207\], Appendix; Acta Phys. Pol. B 31, 1913 (2000) \[hep-ph/0004222\]; cf. also Acta Phys. Pol. B 32, 2961 (2001) \[hep-ph/0108157\]; Acta Phys. Pol. B 33, 2559 (2002) \[hep-ph/0203107\]; hep-ph/0504256.
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# The Color-Superconducting ’t Hooft Interaction
## I History and Motivation
The study of dense (baryon chemical potential $``$ 1.5 GeV) matter has been recently revolutionized by the observation that dense quark matter exhibits color-superconductivity and that the gaps may be of order 100 MeV Bailin and Love (1984); Alford et al. (1998); Rapp et al. (1998). Gaps of this magnitude are large enough to have significant implications for neutron star structure Alford and Reddy (2003), proto-neutron star evolution Carter and Reddy (2000), and neutron star cooling Page et al. (2000).
Since directly utilizing QCD at the relevant densities ($`\mu _{\mathrm{baryon}}1`$ GeV) is so far impossible, the use of effective theories like the Nambu–Jona-Lasinio (NJL) model is common in the study of dense quark matter. In the NJL model, the high-energy degrees of freedom (the gluons) are integrated out and we restrict ourselves to working at energy scales less than the momentum cutoff $`\mathrm{\Lambda }`$ Nambu and Jona-Lasinio (1961); Hatsuda and Kunihiro (1994)
$$_{\mathrm{eff}}=\underset{n}{}\frac{c_n}{\mathrm{\Lambda }^{\mathrm{dim}(𝒪_n)4}}𝒪_n$$
(1)
where $`𝒪_n`$ are operators, $`\mathrm{dim}(𝒪_n)`$ is the dimension of the operator, and $`c_n`$ are dimensionless coupling constants. Because it is impossible to create a model with the same symmetries as QCD with four-fermion operators alone the so-called ’t Hooft term ’t Hooft (1986)
$``$ $`=`$ $`_{\mathrm{kinetic}}+_{\mathrm{four}\mathrm{fermion}}+_{{}_{}{}^{}\mathrm{tHooft}}`$
$`_{{}_{}{}^{}\mathrm{tHooft}}`$ $``$ $`\mathrm{det}_f\left[\overline{q}\left(1+\gamma _5\right)q\right]+\mathrm{det}_f\left[\overline{q}\left(1\gamma _5\right)q\right]`$ (2)
is added where $`q`$ is a quark spinor and $`\mathrm{det}_f`$ is a determinant in flavor space. This term, like QCD, respects chiral symmetry but breaks the $`U(1)_A`$ symmetry. The use of this Lagrangian is the standard approach which has been used to describe dense matter.
Unfortunately, when employed to study quark superconductivity, this standard approach does not employ a manifestly consistent truncation scheme; the quark–anti-quark interaction is treated at the six-fermion level, but the quark-quark interaction is only treated at the four-fermion level. There is no reason to rule out a term of the form (omitting color for clarity)
$`_{\mathrm{CS6}}K_{\mathrm{DIQ}}ϵ_{ijm}ϵ_{k\mathrm{}n}(\overline{q}_i\gamma _5q_j^C)(\overline{q}_k^C\gamma _5q_{\mathrm{}})(\overline{q}_mq_n)`$ (3)
which has the same symmetries as $`_{{}_{}{}^{}\mathrm{tHooft}}`$. This term may have a significant impact on the nature of dense matter Rapp et al. (1998); Schäfer (2002); Buballa (2005). While the effect of the dynamically generated quark masses on the superconducting gaps has been studied Steiner et al. (2002), Eq. 3 implies a modification to the quark masses due to the presence of the gap.
In this article, we study the effect of $`_{\mathrm{CS6}}`$ on the quark masses and on the phase structure of dense matter. We show that sufficiently positive values of $`K_{\mathrm{DIQ}}`$ increase the quark masses and thus favor less-gapped phases, while sufficiently negative values of $`K_{\mathrm{DIQ}}`$ split the magnitude of the up-down and light-strange gaps. The phase structure of dense matter thus depends critically on the sign and magnitude of this unknown parameter.
## II The Model Lagrangian
The Lagrangian is Buballa and Oertel (2002); Gastineau et al. (2002); Steiner et al. (2002)
$``$ $`=`$ $`\overline{q}_{i\alpha }(i/m_{0,ij}\mu _{j\beta }^{i\alpha }\gamma _0)q_{j\beta }+G{\displaystyle \underset{a=0}{\overset{8}{}}}\left(\overline{q}\lambda _f^aq\right)^2`$ (4)
$`+G_{\mathrm{DIQ}}ϵ^{ijm}ϵ^{k\mathrm{}m}ϵ^{\alpha \beta \epsilon }ϵ^{\gamma \delta \epsilon }`$
$`\times \left(\overline{q}_{i\alpha }i\gamma ^5q_{j\beta }^C\right)\left(\overline{q}_{k\gamma }^Ci\gamma ^5q_\mathrm{}\delta \right)`$
$`K\left[\mathrm{det}_f\overline{q}\left(1i\gamma ^5\right)q+\mathrm{det}_f\overline{q}\left(1+i\gamma ^5\right)q\right]`$
$`+K_{\mathrm{DIQ}}ϵ^{ijm}ϵ^{k\mathrm{}n}ϵ^{\alpha \beta \epsilon }ϵ^{\gamma \delta \eta }`$
$`\times (\overline{q}_{i\alpha }i\gamma _5q_{j\beta }^C)(\overline{q}_{k\gamma }^Ci\gamma _5q_\mathrm{}\delta )(\overline{q}_{m\epsilon }q_{n\eta })`$
where flavor is represented by Latin indices, color is represented by Greek indices, and the charge conjugate Dirac spinors are defined by ($`C\gamma ^\mu C=\gamma ^{\mu T}`$ and $`C^T=C`$)
$$q^CC\overline{q}^T\mathrm{and}\overline{q}^C=q^TC.$$
(5)
The four-fermion coupling in the quark–anti-quark channel is denoted $`G`$, the four-fermion coupling in the quark-quark channel is denoted $`G_{\mathrm{DIQ}}`$, and $`m_{0,ij}`$ is the constant current quark mass matrix which is diagonal in flavor.
We utilize the ansätze
$`\overline{q}_1q_2\overline{q}_3q_4`$ $``$ $`\overline{q}_1q_2\overline{q}_3q_4+\overline{q}_1q_2\overline{q}_3q_4\overline{q}_1q_2\overline{q}_3q_4`$
$`\overline{q}_1q_2\overline{q}_3q_4\overline{q}_5q_6`$ $``$ $`\overline{q}_1q_2\overline{q}_3q_4\overline{q}_5q_6+\overline{q}_1q_2\overline{q}_3q_4\overline{q}_5q_6`$ (6)
$`+\overline{q}_1q_2\overline{q}_3q_4\overline{q}_5q_6`$
$`2\overline{q}_1q_2\overline{q}_3q_4\overline{q}_5q_6`$
to obtain the mean-field approximation Bernard et al. (1988). This procedure retains only the lowest order terms in the $`1/N_c`$ expansion Klevansky (1992). Color neutrality is ensured using the procedure from Ref. Steiner et al. (2002). In the mean-field approximation, the Lagrangian is only quadratic in the fermion fields, and the thermodynamical potential can be obtained from the inverse propagator in the standard way Kapusta (1985). The momentum integrals in the gap equations are divergent, and are regulated by a three-momentum cutoff denoted by $`\mathrm{\Lambda }`$. The inverse propagator is numerically diagonalized for each abscissa of the momentum integration to obtain the thermodynamical potential.
For simplicity, we sometimes use the notation $`\mathrm{\Delta }_kϵ^{ijk}q_i^C\gamma ^5q_j`$, so that gaps are denoted with the flavor of quark that is not involved in the pairing e.g. $`\mathrm{\Delta }_{ud}`$ is denoted by $`\mathrm{\Delta }_s`$. Other than $`\overline{q}_iq_i`$ and $`\mathrm{\Delta }_i`$, we assume that all other condensates vanish. This includes the pseudoscalar condensates which which are likely present in dense matter and naturally accompany the Goldstone bosons Kaplan and Reddy (2002). This (non-trivial) complication will be left to later work.
The effects of $`_{\mathrm{CS6}}`$ on the thermodynamical potential can be summarized in three modifications from the standard approach where $`K_{\mathrm{DIQ}}=0`$. These changes are that the values of the gap in the inverse propagator are modified
$`\mathrm{\Delta }_i\mathrm{\Delta }_i\left(1+{\displaystyle \frac{K_{\mathrm{DIQ}}}{N_cG_{\mathrm{DIQ}}}}\overline{q}_iq_i\right),`$ (7)
a new effective mass term (which includes contributions which are not diagonal in flavor) appears
$`{\displaystyle \frac{K_{\mathrm{DIQ}}}{4G_{\mathrm{DIQ}}^2}}\overline{q}_i\mathrm{\Delta }_i\mathrm{\Delta }_jq_j,`$ (8)
which modifies the dynamical mass
$`m_i^{}=m_{i,0}4G\overline{q}q_i+K|ϵ_{ijk}|\overline{q}q_j\overline{q}q_k+{\displaystyle \frac{K_{\mathrm{DIQ}}}{4G_{\mathrm{DIQ}}^2}}\mathrm{\Delta }_i^2`$ (9)
and that there is a new contribution to the vacuum energy
$`\mathrm{\Omega }_{K_{\mathrm{DIQ}}}={\displaystyle \frac{K_{\mathrm{DIQ}}}{2G_{\mathrm{DIQ}}^2}}{\displaystyle \underset{i}{}}\mathrm{\Delta }_i^2\overline{q}_iq_i.`$ (10)
Note that Eq. 8 means that the quark masses are density-dependent if $`K_{\mathrm{DIQ}}0`$ even when the chiral condensates $`\overline{q}q`$ vanish. As has been suggested Rajagopal and Wilczek (2000), this term generates a dynamical quark mass entirely distinct from the typical mechanism of spontaneous chiral symmetry breaking. The quarks obtain a dynamical mass even when the quark condensate is taken to be zero.
One may use Fierz transformations to calculate the coefficient $`K_{\mathrm{DIQ}}`$ from the quark–anti-quark form of the ’t Hooft interaction. One can view this in the following way: For each prospective new term, e.g. the term $`\overline{u}\gamma ^5d_c\overline{d}_c\gamma ^5u\overline{s}s`$, there are six Fierz transformations (for six fermions this is a 35 $`\times `$ 35 matrix instead of the usual 5 $`\times `$ 5 matrix for four fermions) that give this term when applied to the six sets (in flavor space) of terms in the ’t Hooft interaction. One such transformation is
$`\overline{q}_iq_j\overline{q}_kq_{\mathrm{}}\overline{q}_mq_n+\overline{q}_i\gamma ^5q_j\overline{q}_k\gamma ^5q_{\mathrm{}}\overline{q}_mq_n`$ (11)
$`+\overline{q}_iq_j\overline{q}_k\gamma ^5q_{\mathrm{}}\overline{q}_m\gamma ^5q_n+\overline{q}_i\gamma ^5q_j\overline{q}_kq_{\mathrm{}}\overline{q}_m\gamma ^5q_n`$
$`=`$ $`{\displaystyle \frac{1}{2}}(\overline{q}_i\gamma ^5q_k^C\overline{q}_j^C\gamma ^5q_{\mathrm{}}\overline{q}_mq_n+\overline{q}_iq_k^C\overline{q}_j^C\gamma ^5q_{\mathrm{}}\overline{q}_m\gamma ^5q_n`$
$`+\overline{q}_iq_k^C\gamma ^5\overline{q}_j^Cq_{\mathrm{}}\overline{q}_m\gamma ^5q_n+\overline{q}_iq_k^C\overline{q}_j^Cq_{\mathrm{}}\overline{q}_mq_n)`$
$`{\displaystyle \frac{1}{4}}(\overline{q}_i\gamma ^5\sigma ^{\mu \nu }q_k^C\overline{q}_j^C\gamma ^5\sigma _{\mu \nu }q_{\mathrm{}}\overline{q}_mq_n`$
$`+\overline{q}_i\gamma ^5\sigma ^{\mu \nu }q_k^C\overline{q}_j^C\sigma _{\mu \nu }q_{\mathrm{}}\overline{q}_m\gamma ^5q_n)`$
When these transformations are combined to give the coefficient $`K_{\mathrm{DIQ}}`$, the result is zero (see the Appendix). Although terms with different Dirac structure do survive, we do not include these terms since we do not expect the corresponding condensates, e.g. $`\overline{q}\sigma _{\mu \nu }q`$ to be non-zero. This procedure is not the only possible approach for deriving the mean-field Lagrangian (one could also enumerate all possible Wick contractions). Because alternate approaches and/or using terms of higher order in $`1/N_c`$ may modify this result, we cannot conclude necessarily that $`K_{\mathrm{DIQ}}`$ must be zero. Also, terms like Eq. 2 with a Dirac structure $`(\overline{q}q)(\overline{q}\sigma _{\mu \nu }q)(\overline{q}\sigma ^{\mu \nu }q)`$ Shifman et al. (1980) and their corresponding terms in the quark-quark channel may play a role. We leave these considerations to future work.
Our Lagrangian is free to contain any terms which follow the symmetries of the underlying theory. We expect that the coefficient of this term will be “natural”. When the coefficients are expressed in terms of the underlying length scales (in our case, the momentum cutoff $`\mathrm{\Lambda }`$), the coefficients should all be of similar magnitude. We allow the coefficient $`K_{\mathrm{DIQ}}`$ to vary, between the values $`K`$ and $`K`$, which we view to a modest variation as suggested by the constraints of naturalness. A larger variation in $`K_{\mathrm{DIQ}}`$ is not necessarily excluded. We use the values of $`\mathrm{\Lambda }`$ and the current quark masses from Ref. Rehberg et al. (1996), where they are fixed by matching the pion, kaon, and $`\eta ^{}`$ masses in vacuum as well as the pion decay constant. We choose to fix $`G_{\mathrm{DIQ}}\mathrm{\Lambda }^2=1.61`$ to be large enough so that the maximum value of the gaps (when including the quark dynamical mass) as a function of density when $`K_{\mathrm{DIQ}}=0`$ is about 80 MeV, close to the value of about 100 MeV predicted by calculations in perturbative QCD. If the dynamically-generated quark mass is assumed to be zero and the mass-gap equations are ignored, then the maximum value of the gap predicted by this model is about 120 MeV. We leave $`G_{\mathrm{DIQ}}`$ fixed when varying $`K_{\mathrm{DIQ}}`$. One could also allow the coefficient, $`G_{\mathrm{DIQ}}`$ as a function to vary as a function of $`K_{\mathrm{DIQ}}`$ by instead ensuring that the maximum value of one of the three gaps at high densities is constant. We have checked that this alternative does not change our conclusions significantly.
## III Results
Obtaining analytical results is difficult, due to the flavor mixing mass terms from Eq. 8 which make it difficult to directly reduce the inverse propagator (a 36 $`\times `$ 36 matrix) into a block-diagonal form. It is possible, with High-Density Effective Theory Hong (2000), to simplify the Dirac structure, but this would likely result in a 9 $`\times `$ 9 inverse propagator which is also difficult. It is also possible to restore the usual form of the propagator encountered in studies where $`K_{\mathrm{DIQ}}=0`$ by ignoring the terms in Eq. 8 where $`ij`$. In this case, the inverse propagator is worked out in detail in Ref. Rüster et al. (2005).
It is possible to see qualitatively, what the effect of adding a term with $`K_{\mathrm{DIQ}}0`$ might be from Eq. 7 above. When $`K_{\mathrm{DIQ}}>0`$, we expect the gaps decrease as the quark condensate increases, and thus the gap should decrease with increasing mass. However, from Eq. 9 we expect the opposite and we find that it is this effect that dominates the description of the strange quark mass and $`\mathrm{\Delta }_{ud}`$. Further complicating the analysis, Eq. 8 indicates that an increase in $`\mathrm{\Delta }_{us}`$ and $`\mathrm{\Delta }_{ds}`$ will tend to split the mass of the up and down quark, thus possibly weakening $`\mathrm{\Delta }_{ud}`$.
We study charge- and color-neutral, beta-equilibrated, bulk matter at fixed baryon density and a fixed temperature of $`10`$ MeV. We operate at a small but finite temperature in order to alleviate the numerical difficulties of discontinuities in the momentum integral present in the thermodynamical potential. The zero-temperature results will not deviate significantly from our results. We include non-interacting electrons, but we do not include neutrinos. The addition of neutrinos would further split the approximate flavor symmetry between the up and down quarks. Our results will faithfully describe matter in the center of a neutron star containing quarks a minute or later after formation Prakash et al. (1995, 1997).
We note that the effect of $`K_{\mathrm{DIQ}}`$ is small when the quark condensates $`\overline{q}q`$ are taken to be zero. When this is assumed to be true, then the modifications from Eqs. 7 and 8 have no effect, and the gaps are nearly independent of $`K_{\mathrm{DIQ}}`$. We remove this assumption and solve the mass gap equations for the quark condensates in the following.
Figure 1 presents the masses and gaps in the CFL phase at fixed density and temperature as a function of $`K_{\mathrm{DIQ}}`$. Both the quark masses and $`\mathrm{\Delta }_{ud}`$ increase as $`K_{\mathrm{DIQ}}`$ increases. The effect from Eq. 8 causes the $`\mathrm{\Delta }_{us}`$ and $`\mathrm{\Delta }_{ds}`$ gap to decrease when $`\mathrm{\Delta }_{ud}`$ increases. At sufficiently large values of the coupling, the strong increase of the strange quark mass destabilizes the CFL phase. For $`K_{\mathrm{DIQ}}>0.4`$, the gap equations have no solution. If the coupling was verified by some other means to be larger than this critical value, then the CFL phase could not be present at this density. As $`K_{\mathrm{DIQ}}/K1`$, the effects on the masses and gaps tend to be less extreme. The most significant effect is the increasing split between the values of the light-strange gaps, $`\mathrm{\Delta }_{us}`$ and $`\mathrm{\Delta }_{ds}`$, and the light-quark gap $`\mathrm{\Delta }_{ud}`$. One might expect the dependence of the gaps on $`K_{\mathrm{DIQ}}`$ would change the phase structure of matter by shifting the energy density. However, we find that this is not the case and that the energy density is relatively constant as a function of $`K_{\mathrm{DIQ}}`$. Note also that the strange quark mass can change by as much as 50% for different values of $`K_{\mathrm{DIQ}}`$.
Also plotted in Fig. 1 is the parameter $`m_s^2/(\mu \mathrm{\Delta })`$ where $`\mu `$ and $`\mathrm{\Delta }`$ in this context are computed by averaging the quark chemical potentials and gaps over all flavors and colors. This parameter has been demonstrated to be the relevant dimensionless quantity which dictates the phase content of quark matter at high density Alford et al. (2001). Values of $`m_s^2/(\mu \mathrm{\Delta })`$ larger than 2 suggest a transition to a gapless CFL phase Alford et al. (2005), while values larger than 4 suggest a transition to the 2SC phase. The presence of the gapless 2SC phase Shovkovy and Huang (2003) is also probable when $`K_{\mathrm{DIQ}}`$ becomes positive, since the $`ud`$ gap is becoming stronger and the light-strange gaps are weakening. The value of $`m_s^2/(\mu \mathrm{\Delta })`$ is also important in dictating the number and type of Goldstone bosons present in the CFL phase. This parameter is quite flat for small variations in $`K_{\mathrm{DIQ}}`$, but increases or decreases strongly when the absolute magnitude of $`K_{\mathrm{DIQ}}`$ is sufficiently large. When $`K_{\mathrm{DIQ}}/K0.4`$, the value of $`m_s^2/(\mu \mathrm{\Delta })`$ is nearly 2, indicating the disappearance of the CFL phase.
The results for the 2SC phase, where $`\mathrm{\Delta }_{us}=\mathrm{\Delta }_{ds}=0`$, at the same density and temperature are shown in Figure 2. The strange quark mass and gap increase strongly for increasing $`K_{\mathrm{DIQ}}`$ as in the CFL phase, leading to a critical value above which the gap equations have no solution. At this density, the 2SC phase is also not present in matter for $`K_{\mathrm{DIQ}}>0.4`$. The similarity of this critical value of the coupling to the CFL phase in Fig. 1 is a result of the fact that the light-strange gaps and light-quark masses are comparatively small and thus do not strongly perturb $`m_s`$ and $`\mathrm{\Delta }_{ud}`$. The parameter $`m_s^2/(\mu \mathrm{\Delta })`$ is slightly modified from the CFL case and indicates that the 2SC phase is also likely to be unstable for large negative values of $`K_{\mathrm{DIQ}}/K`$.
For comparison, Fig. 3 presents the quark masses and gaps in the CFL phase at a lower density, 1.2 fm<sup>-3</sup>. The results are not much different from Fig. 1. The quark gaps have decreased slightly and the quark masses are slightly larger. The major distinction is that the critical value of $`K_{\mathrm{DIQ}}`$ above which the CFL phase does not appear has been lowered from 0.4 to less than 0.3.
The implication of the shift in the critical value of $`K_{\mathrm{DIQ}}`$ is that the critical density for the onset of the CFL phase is drastically affected by a positive value of $`K_{\mathrm{DIQ}}`$. This is demonstrated in Fig. 4, where the smallest quark gap vanishes at about 1.0 fm<sup>-3</sup> when $`K_{\mathrm{DIQ}}=0`$, and at about 1.42 fm<sup>-3</sup> when $`K_{\mathrm{DIQ}}/K=0.3`$. On the other hand, because the gaps are not as strongly modified when $`K_{\mathrm{DIQ}}`$ is negative, the critical density when $`K_{\mathrm{DIQ}}/K=0.5`$ is almost unchanged, moving down only to 0.96 fm<sup>-3</sup>. In this figure, the $`\mathrm{\Delta }_{ds}`$ gap does not vanish completely to zero because the solution of the gap equations becomes difficult when the gaps are small. Note again that an increase the parameter $`m_s^2/(\mu \mathrm{\Delta })`$ indicates, to some extent, the disappearance of the CFL phase as the density decreases.
## IV Discussion - Neutron Stars with $`K_{DIQ}0`$
These results may have several implications for neutron star structure and evolution.
Neutron-star masses and radii: We have solved the Tolman-Oppenheimer-Volkov equations for neutron star structure under the assumption that there is no mixed phase (i.e. the surface tension is very large so that mixed phases are suppressed). The results for $`K_{\mathrm{DIQ}}=0`$ and $`K_{\mathrm{DIQ}}/K=0.3`$ are given in Fig. 5. We find that neutron star masses and radii are not very sensitive to $`K_{\mathrm{DIQ}}`$ for the model that we have chosen. As mentioned above, a positive value of $`K_{\mathrm{DIQ}}`$ tends to increase the critical density for the appearance of the CFL phase, thus stiffening the equation of state and slightly increasing the maximum mass from 1.83 to 1.9 $`M_{}`$. The small magnitude of this effect is partially due to the fact that, for $`K_{DIQ}=0`$, the critical density for the appearance of quark matter is 1.0 fm<sup>-3</sup>, while the central density of the maximum mass neutron star is only 1.45 fm<sup>-3</sup>. There is not much quark matter to begin with, so therefore the effect of $`K_{\mathrm{DIQ}}`$ is limited. In regards to the phase content of matter in the neutron star, the larger value of $`K_{\mathrm{DIQ}}`$ nearly pushes the CFL phase out of the neutron star entirely, and the center of the neutron star is dominated by the 2SC phase instead. These effects will be larger if the diquark coupling is increased and may be modified by the presence of a mixed phase. Recent calculations of the surface tension suggest that it is small, and thus a mixed phase may be present Reddy and Rupak (2005). It would be interesting to examine the effect of $`K_{\mathrm{DIQ}}`$ on this surface tension.
Neutron-star cooling: Neutron stars are sensitive to the difference between the CFL and 2SC phases since the former is likely to contain no Rajagopal and Wilczek (2001) (or very few) electrons. Neutron star cooling is affected by the presence or absence of electrons because the specific heat of matter is dominated by the electrons when they are present. In addition, the specific heat contribution from the light quarks is proportional to $`\mathrm{exp}(\mathrm{\Delta }_\mathrm{}s/T)`$, where $`\mathrm{}=u\mathrm{or}d`$ which is exponentially small in the CFL phase and of order unity in the 2SC phase. Also, the splitting of the gaps at $`K_{\mathrm{DIQ}}/K=1`$, will enforce two critical temperatures for quark matter in the CFL phase. The so-called “gapless” phases are dependent upon the strange quark mass which is strongly affected by $`K_{\mathrm{DIQ}}`$, especially when it is greater than zero. These gapless phases, because of the nature of the quark dispersion relations, have unusual transport properties that also have implications for neutron star cooling Shovkovy and Huang (2003); Huang and Shovkovy (2004); Alford et al. (2005).
Proto-neutron star evolution: Ref. Carter and Reddy (2000) pointed out that the proto-neutron star neutrino signal may be increased noticeably by the enhanced cooling that is present when the core temperature falls below the critical temperature. The presence of $`K_{\mathrm{DIQ}}`$ implies that this enhanced cooling may occur in (at least) two separate stages, as the critical temperatures corresponding to the $`\mathrm{\Delta }_{ud}`$ and the light-strange gaps are surpassed.
To the extent that the presence of a large ($`100`$ MeV) gap affects observations of neutron stars, the presence of a color-superconducting six-fermion interaction also has an impact on neutron star observations. It would be interesting to compare these results with those from gapless CFL phases. Because of the uncertainty in the value of the $`K_{\mathrm{DIQ}}`$ coupling (in addition to the other uncertainties already present), it will be difficult to settle the questions of the nature of dense matter and the question of the properties of neutron stars containing deconfined quark matter until theoretical progress in QCD or astrophysical observations can more completely settle the issue.
The author would like to thank T. Bhattacharaya, T. Bürvenich, M. Prakash, S. Reddy, G. Rupak, and T. Schäfer for discussions. This work has been supported by the DOE under grant number DOE/W-7405-ENG-36.
## Appendix A - Fierz Transformations
Traditionally, the Fierz transformation is defined as the matrix, $`C`$, which obeys the relation
$`\left[\overline{\psi }_i\mathrm{\Gamma }_{a,ij}\psi _j\overline{\psi }_k\mathrm{\Gamma }_k\mathrm{}^a\psi _{\mathrm{}}\right]={\displaystyle \underset{b}{}}C_{a,b}\overline{\psi }_i\mathrm{\Gamma }_{b,i\mathrm{}}\psi _{\mathrm{}}\overline{\psi }_k\mathrm{\Gamma }_{kj}^b\psi _j`$ (12)
where $`\mathrm{\Gamma }_i[11,\gamma ^5,\gamma _\mu ,\gamma ^5\gamma _\mu ,\sigma _{\mu \nu }]`$ and $`\mathrm{\Gamma }^i[11,\gamma ^5,\gamma ^\mu ,\gamma ^5\gamma ^\mu ,\sigma ^{\mu \nu }]`$ for $`i=a,b`$. Because the Fierz transformation is nothing other than a set of equalities, the two four-fermion interactions on both sides are necessarily equivalent. However, in the mean-field approximation, these two forms lead to different thermodynamic potentials. For simplicity, the transformation above will be denoted $`(ijk\mathrm{}i\mathrm{}kj)`$.
One may also perform a four-fermion Fierz transformation in the quark-quark channel, namely, $`(ijk\mathrm{}ikj\mathrm{})`$ (see the review in Ref. Nieves and Pal (2004)). This transformation operates over a different “basis” of combinations of Dirac matrices
$`\left[\overline{\psi }_i\mathrm{\Gamma }_{a,ij}\psi _j\overline{\psi }_k\mathrm{\Gamma }_k\mathrm{}^a\psi _{\mathrm{}}\right]=`$
$`{\displaystyle \underset{b}{}}C_{a,b}^{}\overline{\psi }_i\mathrm{\Gamma }_{b,ik}^{}\psi _k^C\overline{\psi }_j^C\mathrm{\Gamma }_j\mathrm{}^{,b}\psi _{\mathrm{}}.`$ (13)
where $`\mathrm{\Gamma }_i^{}[\gamma _5,11,\gamma ^5\gamma _\mu ,\gamma _\mu ,\sigma _{\mu \nu }]`$ and $`\mathrm{\Gamma }^{i,}[\gamma ^5,11,\gamma ^5\gamma ^\mu ,\gamma ^\mu ,\sigma ^{\mu \nu }]`$ (remember that $`\overline{\psi }^C\gamma ^5\psi `$ is a Lorentz scalar). We can simplify the notation for the basis by using the notation of a direct product: $`\gamma ^5\gamma ^5,1111,\gamma ^5\gamma _\mu \gamma ^5\gamma _\mu ,\gamma _\mu \gamma _\mu ,\sigma _{\mu \nu }\sigma ^{\mu \nu }`$, or more simply, $`\mathrm{SS},\mathrm{PP},\mathrm{VV},\mathrm{AA},\mathrm{TT}`$.
In a similar notation the 35-element basis for the six-fermion Fierz transformations is
$`\mathrm{SSS},\mathrm{SPP},\mathrm{PSP},\mathrm{PPS},\mathrm{SVV},\mathrm{VSV},\mathrm{VVS},`$
$`\mathrm{SAA},\mathrm{ASA},\mathrm{AAS},\mathrm{STT},\mathrm{TST},\mathrm{TTS},`$
$`\mathrm{PVA},\mathrm{PAV},\mathrm{VPA},\mathrm{VAP},\mathrm{APV},\mathrm{AVP},`$
$`\mathrm{TVV},\mathrm{VTV},\mathrm{VVT},\mathrm{TAA},\mathrm{ATA},\mathrm{AAT},\mathrm{TTT},`$
$`\mathrm{VAQ},\mathrm{VQA},\mathrm{AVQ},\mathrm{AQV},\mathrm{QVA},\mathrm{QAV},`$
$`\mathrm{PTQ},\mathrm{TPQ},\mathrm{QTP}`$ (14)
where $`Q`$ denotes a “pseudo-tensor” combination, $`\gamma ^5\sigma ^{\mu \nu }`$. (The $`\mathrm{Q}`$ terms are used as an alternative to the formulation in terms of objects of the form $`\epsilon _{\kappa \lambda \mu \nu }\sigma ^{\mu \nu }`$ Maruhn et al. (2001).) All other combinations can be written as a linear combination of these 35 basis elements. In order to distinguish 4- and 6-fermion transformations, we will use a subscript, i.e. $`_4(ijk\mathrm{}mni\mathrm{}kjmn)`$ is really a four-fermion transformation since the fields with indices $`m`$ and $`n`$ are not participating in the transformation.
The Dirac scalar terms in the ’t Hooft interaction are
$`\overline{u}u\overline{d}d\overline{s}s+\overline{u}s\overline{d}u\overline{s}d+\overline{u}d\overline{d}s\overline{s}u`$ (15)
$`\overline{u}d\overline{d}u\overline{s}s\overline{u}s\overline{d}d\overline{s}u\overline{u}u\overline{d}s\overline{s}d`$ $`.`$
plus the corresponding terms created by adding an even number of $`\gamma ^5`$ matrices.
As a demonstration, we examine the coefficient of the term $`\overline{u}\gamma ^5d^C\overline{u}^C\gamma ^5d\overline{s}s`$. The various contributions to this coefficient are
$`+_4(ijklmnikjlmn)`$ $`[\overline{u}u\overline{d}d\overline{s}s+\overline{u}u\overline{d}\gamma ^5d\overline{s}\gamma ^5s`$
$`+\overline{u}\gamma ^5u\overline{d}d\overline{s}\gamma ^5s+\overline{u}\gamma ^5u\overline{d}\gamma ^5d\overline{s}s]`$
$`+_6(ijklmniklnmj)`$ $`[\overline{u}s\overline{d}u\overline{s}d+\overline{u}s\overline{d}\gamma ^5u\overline{s}\gamma ^5d`$
$`+\overline{u}\gamma ^5s\overline{d}u\overline{s}\gamma ^5d+\overline{u}\gamma ^5s\overline{d}\gamma ^5u\overline{s}d]`$
$`+_6(ijklmniknjml)`$ $`[\overline{u}d\overline{d}s\overline{s}u+\overline{u}d\overline{d}\gamma ^5s\overline{s}\gamma ^5u`$
$`+\overline{u}\gamma ^5d\overline{d}s\overline{s}\gamma ^5u+\overline{u}\gamma ^5d\overline{d}\gamma ^5s\overline{s}u]`$
$`_4(ijklmnikljmn)`$ $`[\overline{u}d\overline{d}u\overline{s}s+\overline{u}d\overline{d}\gamma ^5u\overline{s}\gamma ^5s`$
$`+\overline{u}\gamma ^5d\overline{d}u\overline{s}\gamma ^5s+\overline{u}\gamma ^5d\overline{d}\gamma ^5u\overline{s}s]`$
$`_6(ijklmniknlmj)`$ $`[\overline{u}s\overline{d}d\overline{s}u+\overline{u}s\overline{d}\gamma ^5d\overline{s}\gamma ^5u`$
$`+\overline{u}\gamma ^5s\overline{d}d\overline{s}\gamma ^5u+\overline{u}\gamma ^5s\overline{d}\gamma ^5d\overline{s}u]`$
$`_6(ijklmnikjnml)`$ $`[\overline{u}u\overline{d}s\overline{s}d+\overline{u}u\overline{d}\gamma ^5s\overline{s}\gamma ^5d`$
$`+\overline{u}\gamma ^5u\overline{d}s\overline{s}\gamma ^5d+\overline{u}\gamma ^5u\overline{d}\gamma ^5s\overline{s}d]`$
Note that result of the first of these six transformations is given as Eq. 11 in the text. The coefficients of the desired term, $`\overline{u}\gamma ^5d^C\overline{u}^C\gamma ^5d\overline{s}s`$, from each of the six transformations (together with an appropriate factor of -1 for odd fermionic permutations) are
$$\frac{1}{2},\frac{1}{4},\frac{1}{4},\frac{1}{2},\frac{1}{4},\frac{1}{4}$$
(16)
and the sum is zero.
Because the coefficient is zero, we need not consider the Fierz transformations in the SU(3) (color or flavor) spaces. However, since the result for six-fermion transformations in SU(3) is not present in the literature, we give the 15-element basis for computing the transformations (this enlarged basis from that presented in Ref. Dmitrasinovic (2001) is necessary to perform the transformations in the quark–quark channel)
$`111111,\lambda ^a\lambda _a11,`$
$`11\lambda ^a\lambda _a,\lambda _a11\lambda ^a,`$
$`d_{abc}\lambda ^a\lambda ^b\lambda ^c,if_{abc}\lambda ^a\lambda ^b\lambda ^c,`$
$`\lambda (A)^a\lambda (A)_a11,11\lambda (A)^a\lambda (A)_a,`$
$`\lambda (A)_a11\lambda (A)^a,`$
$`\lambda (A)_a\lambda (S)_b\lambda (A)^a\lambda (S)^b,`$
$`\lambda (A)_a\lambda (S)_b\lambda (S)^b\lambda (A)^a,`$
$`\lambda (S)^b\lambda (A)_a\lambda (S)_b\lambda (A)^a,`$
$`\lambda (A)^a\lambda (A)_a\lambda (S)_b\lambda (S)^b,`$
$`\lambda (A)^a\lambda (S)^b\lambda (A)_a\lambda (S)_b,`$
$`\lambda (S)^b\lambda (A)^a\lambda (A)_a\lambda (S)_b`$ (17)
where the occurrence of $`\lambda `$ indicates an implicit sum over all 8 SU(3) matrices, $`\lambda (A)`$ restricts the sum to only the three anti-symmetric $`\lambda `$ matrices, and the implicit sum for $`\lambda (S)`$ is over the six symmetric $`\lambda `$ matrices. The full results are available from the author.
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# Effects of clumping on temperature I: externally heated clouds
## 1 Introduction
An understanding of the star formation process requires an understanding of the underlying density distribution in star forming regions. In a direct sense, knowledge of the density distribution can help distinguish between different potential dynamical scenarios for energy injection, collapse, fragmentation, and outflow. In a more indirect sense, the density distribution significantly affects our ability to infer source properties through its influence on thermal balance, chemistry, local emission, and the processing of radiation (radiative transfer) between the emitting region and the observer.
The density distribution is important for the dust as well as the gas. In particular, the dust forms the dominant source of opacity to visible and infrared (IR) radiation, and is the dominant source of IR continuum radiation. Perhaps more importantly, the dust dominates thermal balance by direct interaction with the radiation field (van de Hulst 1949), and through collisions with the gas (Greenberg 1971; Goldreich & Kwan 1974). As the problems of thermal balance and radiative transfer are non-local, non-linear feedback problems, comparison of detailed models with observations remain the best choice of reliably inferring the source properties.
Source geometry is a significant problem in modeling star-forming regions. While it is normal to assume some geometric symmetry, such restrictions are generally not realistic. In particular, a wealth of observations (e.g. Migenes et al. 1989; Dickman et al. 1990; Falgarone et al. 1991; Cesaroni et al. 1991; Marscher et al. 1993; Zhou et al. 1994; Plume et al. 1997; Shepherd et al. 1997) show that and fragmentation (see also Goldsmith 1996; Tauber 1996 and references therein). This is supported by dynamical models (e.g. Truelove et al. 1998; Marinho & Lepine 2000; Klapp & Sigalotti 1998; Klessen 1997) which naturally produce clumpy and fragmented structures.
Previous work on radiative transfer in dusty, clumpy environments has been undertaken (e.g. Hegmann & Kegel 2003; Witt & Gordon 1996; Varosi & Dwek 1999; Boissé 1990). However, much of the previous work has concentrated on scattering (e.g. with application to reflection nebulae), the detailed methods of solution, and/or some of the fundamental results such as ability of radiation to penetrate to apparently high optical depths. In this paper we use a 3-D monte carlo radiative transfer model to extend the previous work to a wider and different range of physical conditions. In particular, we construct a large grid of models in an effort to better delineate and understand the effects of clumpy media on radiative transfer. By controlling and varying the parameterization of the source, we attempt to disentangle some of the underlying physical causes of these effects.
In section two we describe the model. We discuss the general effects of clumping in section three. In section four, we introduce the “fraction of open sky” (FOS), and discuss the effects of number density, FOS, optical depth, filling factor, and clump size on the dust temperature distribution. We discuss the effects of shadowing in section five. Finally, we draw conclusions in section six.
## 2 Model
### 2.1 Monte carlo model and invariant parameters
We have constructed detailed, self-consistent, three-dimensional radiative transfer models through dust. The model utilizes a monte-carlo approach, combined with an approximate lambda iteration to ensure true convergence even at high optical depths. The model has been tested against existing 1-D (Egan, Leung, & Spagna 1988) and 2-D (Spagna, Leung, & Egan 1991) codes, and in modeling a 3-D source (Doty et al. 2005) with good success. Since the present study concentrates upon opaque molecular clouds where far-infrared radiaton should dominate both for external heating and emission, we ignore the effects of scattering. Test cases in both one- and multiple-dimensions show that scattering plays only a very small role on the temperature distribution within the majority of the sources.
Based upon the input parameters discussed below, we solve for the dust temperature and radiation field at each point in the model cloud. The computational volume is taken to be cubical of size 1 pc. Each model utilized an $`81\times 81\times 81`$ cell grid, yielding a typical resolution of $`3\times 10^{16}`$ cm. The region is taken to be a two-phase medium consisting of high density clumps, and a low density inter-clump medium. The clump/inter-clump density ratio is taken to be $`n_{\mathrm{clump}}/n_{\mathrm{inter}\mathrm{clump}}=100`$ from observations (e.g. Bergin et al. 1996). The external radiation field is taken from interstellar radiation field (ISRF) compiled by Mathis, Mezger, & Panagia (1983). Finally, we adopt the dust opacities in column 5 of Table 1 of Ossenkopf & Henning (1994), which have been successful at fitting observations of both high-mass (e.g., van der Tak et al. 1999, 2000) and low-mass (e.g., Evans et al. 2001) regions of star formation.
### 2.2 Model parameters
We adopt a uniformly distributed interclump medium interspersed with higher density clumps having $`n_{\mathrm{clump}}/n_{\mathrm{inter}\mathrm{clump}}=100`$. The clumps are randomly distributed within the computational volume. The number of clumps, and the densities of the clumps and interclump medium depend upon the filling factor ($`f`$), the clump radius ($`r_c`$), and the face-averaged optical depth ($`\overline{\tau }`$).
The filling factor specifies the fraction of the volume at high density. It is given by $`fV_{\mathrm{clumps}}/V_{\mathrm{total}}`$. We adopt filling factors of $`f=0.01`$, $`0.1`$, and $`0.3`$ in accord with observations (e.g. Snell et al. 1984; Bergin et al. 1996; Carr 1987). When all other parameters are kept fixed, a higher $`f`$ corresponds to a larger number of clumps, and a more nearly continuous dust distribution. Furthermore, due to optical depth constraint (see $`\overline{\tau }`$ below), a larger $`f`$ also corresponds to clumps and interclump medium of lower density.
We make the simplifying assumption that all non-overlapping clumps are spheres of radius $`r_c`$. We choose $`r_c=0.025`$ pc, $`0.05`$ pc, and $`0.1`$ pc in accord with observations (e.g. Carr 1987; Howe et al. 1993). Again, the constraint on the optical depth implies that larger clump radii yield smaller densities within the clumps.
The dust number densities are normalized by averaging the optical depth over one entire ($`81\times 81`$) face of the computational cube. We call this the face-averaged optical depth, and denote it by $`\overline{\tau }`$. This was done to simulate the optical depth / column density that might be inferred by a very large beam, although it has the same effect as normalizing to the total cloud mass. The face-averaged optical depth is taken to be $`\overline{\tau }=10`$, and $`100`$, in keeping with observations from extinction studies (Lada, Alves, & Lada 1999).
Finally, the models are specified by the strength of the internal radiation field, specified by the luminosity of the central source, $`L_{}`$. We have constructed a grid having $`L_{}=0`$L, $`3`$L, and $`300`$L to represent a starless core, a low-luminosity central object, and a high-luminosity central object, respectively. However, in this paper, we restrict our report to the starless cores only ($`L_{}=0`$). The others will be presented in a forthcoming paper.
The ranges of parameters specified above led to a grid of 54 models. Each model is numbered by a four-digit integer $`I_LI_fI_{rc}I_{\overline{\tau }}`$. The integers, and their corresponding values are given in Table 1 for reference. As an example, model 1221 has $`L_{}=0`$, $`f=0.01`$, $`r_c=0.025`$ pc, and $`\overline{\tau }=10`$.
### 2.3 Invariance with clump and photon randomization
One concern with “random” clump distribution and monte-carlo simulation is the reproducibility of results with different realizations of the clump distribution (i.e. different initial seeds in the clump generation), and different photon ray paths. As a test, we have considered 9 different initial seeds for the clump distributions and photon paths respectively.
To quantify the temperature distribution, we calculate a spherical average temperature, given by
$$T(r)=\underset{i}{\overset{N}{}}T_in_i/\underset{i}{\overset{N}{}}n_i.$$
(1)
Here $`N`$ is the number of cells a distance between $`r`$ and $`r+\mathrm{\Delta }r`$ from the center, where $`\mathrm{\Delta }r`$ is the size of a single cell, and $`T_i`$ and $`n_i`$ are the temperature and density of the cells respectively. In keeping with the viewpoint that the high density clumps mostly affect the local radiation field, while the the low density medium best probes the radiation field, the average temperature is taken only over the low-density cells.
The differences in $`T(r)`$ between different clump distribution realizations are always much less than 1K, and are on average less than 0.1K. Based upon this (and direct comparison of later analyses) the specific realization of the density distribution does not affect our conclusions.
Likewise, we have modified the number and distribution of incident photons to test for sufficient monte carlo coverage. The effect of number of photons on for $`T(r)`$ is shown in Fig. 1. To be conservative, we adopt 100 photons per cell face, at which point differences are less than 0.05K (1 per cent).
Similarly, varying the random distribution of photon paths causes deviations of $`<0.1`$K, confirming that the coverage of the cells by ray-paths is sufficient to concentrate on the consequences of clumping.
## 3 Clumping: general
In this section we briefly discuss the general effects of clumping on the dust temperature distribution. For clarity, we directly compare a clumpy, externally heated ($`L_{}=0`$), low-density ($`f=0.1`$), opaque ($`\overline{\tau }=10`$), model having small clumps ($`r_c=0.05`$ pc) to one with a uniform density distribution and the same optical depth and external heat source, so we can directly measure the effects of clumping. We have done similar comparisons for all other models, with similar results.
The temperature distributions along the principal axes for a representative clumpy and equivalent homogeneous model are shown in Fig. 2. To quantify the differences between the temperature distributions, in Fig. 3 we plot the fractional difference in $`T(r)`$ between these two representative models. From these two figures, it is immediately obvious that the inclusion of clumps – even for the same $`\overline{\tau }`$ – changes the temperature structure significantly. In particular, the clumpy model experiences higher temperatures by up to 50 per cent toward the center, and up to $`65`$ per cent at intermediate radii. As a result, we conclude that clumping itself affects the temperature distribution, even when the average source mass or column density is held constant.
This result confirms the previous finding of others (e.g. Hegmann & Kegel 2003; Witt & Gordon 1996; Varosi & Dwek 1999) that the effective optical depth in a clumpy medium is less than the homogeneous value. It also suggests that it is important to understand the way in which the parameterization of the clumpy density distribution can affect the dust temperature. We address these individually below.
## 4 Clumping: effects of parameters
### 4.1 Number density
In Fig. 4 we plot the temperature (solid lines, left hand scale) and dust number density (dotted lines, right hand scale) for cuts along the x-, y-, and z-axes. The clumps are signified by the higher density regions. As can be seen, the clumps are resolved, and appear to be of different sizes as the axes do not penetrate all clumps along a diameter. The “wiggles” and $`0.10.2`$ K deviations in the temperature distribution are not simply indicative of the uncertainties in the monte carlo calculation (see previous results, and Fig. 2). Instead, the majority are dominated by local differences in radiation field due to different amounts of blocking of radiation by the surrounding clumps (see Sect. 5 for more discussion).
As can be seen in Fig. 4, the dust temperature is lower within the high density clumps than it is within the lower-density interclump medium. Since the total emission by grains of absorption efficiency $`Q_{\mathrm{abs}}\pi a^2\nu ^\beta `$ goes as $`T^{4+\beta }`$, one might expect a factor of 100 increase in density to lead to a decrease in temperature by a factor $`100^{1/(4+\beta )}`$. Given that $`\beta 1.8`$ in the FIR for the adopted dust model, this is a decrease by a factor of $`2.2`$. For comparison, the average decrease within the clumps is a factor of $`1.5`$ with a greatest decrease a factor of $`1.6`$. This is due to the fact that the clump centers can be warmed by absorbing the radiation emitted from the clump edges, in effect trapping radiation within the clumps (similar to the opaque centers of centrally heated envelopes, Doty & Leung 1994).
Of even further interest is the fact that while the temperature decreases coincide with the locations of the clumps, the shape of the temperature profiles do not match the steepness of the clump/interclump interfaces. This further suggests that radiative transfer effects play a role.
### 4.2 Effective optical depth
Previous authors (e.g., Hegmann & Kegel 2003; Witt & Gordon 1996; Varosi & Dwek 1999) identified the effective optical depth ($`\mathrm{ln}[L_{\mathrm{AT}}/L_{}]`$), where $`L_{\mathrm{AT}}`$ is the attenuated flux integrated over all solid angles and $`L_{}`$ is the unattenuated flux integrated over all solid angles or equivalently the average attenuation ($`e^\tau e^\tau 𝑑\mathrm{\Omega }/𝑑\mathrm{\Omega }`$) as a key measure of the ability of radiation to penetrate. This is reasonable, as the fewer photons penetrating leads to cooler dust, which is confirmed in Fig. 5, where we plot the dust temperature at three different positions along the x-axis as a function of $`e^\tau `$ for model 1221.
One question is left unanswered – the physical process driving the angle-averaged attenuation must be identified. There are two possibilities: (1) the ability of radiation to penetrate is dominated by general attenuation by multiple clumps along a given line of sight and attenuation by the low-density interclump medium; or (2) the radiation streams primarily through the “holes” between clumps. We consider these two possibilities in the following subsection.
### 4.3 Optical depth and fraction of open sky
#### 4.3.1 Motivation and definition of $`\tau `$ and FOS
In order to answer the question of whether and how much geometry matters in determining the local temperature/radiation field, we consider two limiting cases. These are a measure with limited geometry information, the angle-averaged optical depth ($`\tau `$), and a measure which is dominated by geometry information, the fraction of open sky (FOS). We consider these seperately below.
If the radiation generally diffuses and is attenuated by multiple clumps or the low-density interclump medium, we might expect to find a dependence of source properties on the optical depth averaged over all directions. As a result, we define the angle-averaged optical depth to be
$$\tau \tau (\mathrm{\Omega })𝑑\mathrm{\Omega }/𝑑\mathrm{\Omega }.$$
(2)
On the other hand, if the radiation penetrates mainly by streaming through holes between clumps, then the source properties should depend more significantly on the fraction of the sky that is uncovered by the clumps. Consequently, we also define the fraction of open sky, FOS, to be
$$\mathrm{FOS}q𝑑\mathrm{\Omega }/𝑑\mathrm{\Omega }.$$
(3)
Here $`q=1`$ if there is no clump along the given line of sight (i.e. if the sky is “open” in that direction), while $`q=0`$ if there is a clump along the given line of sight (i.e. if the sky is “closed” in that direction). Consequently, FOS=1 implies an open sky with no clumps, while FOS=0 corresponds to a closed sky where all lines of sight are blocked by clumps.
As a first test of the effect of FOS on temperature, in Fig. 6 we plot the temperature (solid lines, left-hand scale) and FOS (dotted lines, right-hand scale) for cuts along the x-, y-, and z-axes, similar to Fig. 4. There is a strong correlation of dust temperature with FOS. Interestingly, the shape of the temperature distribution correlates with FOS to a much higher degree than it anti-correlates with the dust number density. This strongly suggests that FOS plays a key role, which we test in the following discussion.
#### 4.3.2 Effect of $`\tau `$ on temperature
To test the role of angle-averaged optical depth, we consider the dependence of temperature on $`\tau `$ for a fixed FOS. In this way, we can isolate the effects of $`\tau `$ from FOS. The FOS value here was chosen to maximize the number of data points to provide for the best possible statistics. The resulting dependence of temperature on $`\tau `$ for three different radial positions is shown in Fig. 7. We note that the chosen FOS range is small to make the result independent of FOS, though moderate increases in the range lead to similar results.
There appears to be little correlation of temperature with $`\tau `$ at a given position. We quantify this by fitting the data distribution at each radial distance with a best-fit straight line. To understand the significance of the slope / correlation, we define a correlation parameter,
$$\theta (\tau )m(\tau )/\sigma _{m(\tau )}.$$
(4)
Here, $`m(\tau )`$ is the slope of the best-fit line relating the temperature and $`\tau `$, and $`\sigma _{m(\tau )}`$ is the uncertainty in that slope. The results are shown in Fig. 8. We have also included dashed lines at $`m=\pm 3\sigma _{m(\tau )}`$ as a guide.
To best understand the utility of this comparison, consider a case for which the temperature and $`\tau `$ are uncorrelated. In this case, the fit between them should yield a zero slope, and thus $`\theta (\tau )=0`$. Likewise, a significant correlation should yield $`\theta >3`$ (i.e., a $`3\sigma `$ result). The results in Fig. 8 are much more consistent with $`\theta (\tau )=0`$ throughout much of the cloud, and only reach $`\theta (\tau )=3`$ for $`r>0.7R_{\mathrm{out}}`$. Consequently, we conclude that there may exist only a weak correlation between the dust temperature and $`\tau `$.
#### 4.3.3 Effect of FOS on temperature
In a similar manner to $`\tau `$, we consider the effect of FOS on temperature. In this case, we keep $`\tau `$ fixed. In analogy with before, $`\tau `$ was chosen to maximize the number of data points to ensure the best possible statistics. The resulting dependence of temperature on FOS for three different radial positions is shown in Fig. 9.
Inspection of Fig. 9 suggests that there exists some correlation of temperature with FOS. In a similar manner to $`\tau `$, we quantify this by fitting the data distribution at each radial distance with a best-fit straight line. We again define a correlation parameter,
$$\theta (\mathrm{FOS})m(\mathrm{FOS})/\sigma _{m(\mathrm{FOS})}.$$
(5)
Here, $`m(\mathrm{FOS})`$ is the slope of the best-fit line relating the temperature and FOS, and $`\sigma _{m(\mathrm{FOS})}`$ is the uncertainty in that slope. The results are show in Fig. 10. We have also included dashed lines at $`\pm 3\sigma _{m(\mathrm{FOS})}`$ as a guide.
The results in Fig. 10 are generally not consistent with $`\theta (\mathrm{FOS})=0`$, and lie at or above $`\theta (\mathrm{FOS})=3`$ for a good deal of the cloud ($`r>0.25R_{\mathrm{out}}`$). Consequently, we can say that there exists a significant correlation between the dust temperature and FOS.
Finally, it is encouraging to note that the correlation between temperature and FOS is positive. This is expected, since a higher FOS leads to more direct heating by the external radiation field, and is evidenced by $`m(\mathrm{FOS})>0`$.
#### 4.3.4 Extension and interpretation of $`\tau `$, FOS results
We can extend the discussion of FOS and $`\tau `$ to the remainder of our family of models. The results are summarized in Table 2. In this table, we specify the fractional radial position from cloud center as $`xr/R_{\mathrm{out}}`$.
The results in Table 2 require an understanding of the range over which FOS or $`\tau `$ can have a significant effect on the temperature distribution. In the table, $`x_{\mathrm{max}}`$ is the point beyond which, FOS does not play a significant role on the temperature distribution. We find (see Section 4.5) this occurs for $`\mathrm{FOS}>0.60.8`$.
To understand the origin and location of this region, consider a point a distance $`r`$ from the center of the region of radius $`R_{\mathrm{out}}`$. The number of clumps exterior to $`r`$ is $`N=f(R_{\mathrm{out}}/r)^3(1x)`$. If these clumps are distributed about this point on the outward facing half of an equal volume sphere of radius $`r_{\mathrm{eff}}=[4R_{\mathrm{out}}^3(1x)/(4\pi /3)]^{1/3}`$, then the fraction of open sky in the outward direction is $`\mathrm{FOS}^+=1Nr_c^2/2r_{\mathrm{eff}}^2`$. We plot the resulting estimate of FOS<sup>+</sup> as a function of position in Fig. 11. As can be seen, all models with a low filling factor (x1xx) have FOS$`{}_{}{}^{+}>0.8`$, as does model (x23x). For these models, FOS is not expected to play a significant role in determining the temperature distribution. As a result, the radius up to which FOS should play a role ($`x_{\mathrm{max}}`$) becomes zero for these models in Table 2. For other models, $`0<x_{\mathrm{max}}<1`$.
To understand the comparative importance of FOS vs. $`\tau `$, we consider the average deviation of $`\theta `$ from zero for $`x<x_{\mathrm{max}}`$. This is reported in the last column of Table 2. For models in which $`x_{\mathrm{max}}=0`$, such an average is not possible, and the result is noted as “n/a”. For the case where only $`\theta (\mathrm{FOS})>3`$ for $`x<x_{\mathrm{max}}`$, the only meaningful correlation of temperature is with FOS, and so the result is noted as $`\mathrm{}`$. In all other cases, the last column is set to $`\overline{\theta (\mathrm{FOS})}/\overline{\theta (\tau )}`$, where the bar signifies an average over all positions for which $`\theta >3`$ and $`x<x_{\mathrm{max}}`$. As can be seen by the last column of Table 2 the effect of FOS exceeds the effects of $`\tau `$ in determining the temperature.
Finally, one may consider the size of the region over which FOS and $`\tau `$ are important. To do this, in Table 2 we specify $`x_{\mathrm{cut}}`$, defined to be the position beyond which $`\theta >3`$. Generally $`x_{\mathrm{cut}}(\tau )>x_{\mathrm{cut}}(\mathrm{FOS})`$. At low optical depths, this is true whenever $`x_{\mathrm{max}}0`$. For this set, the mean values for $`x_{\mathrm{cut}}(\tau )0.7`$, and $`x_{\mathrm{cut}}(\mathrm{FOS})0.4`$. This suggests that not only is FOS (and thus photon streaming) more important than $`\tau `$, but that it is also important deeper into the cloud than $`\tau `$. This is not a suprise, as openings are a significantly more efficient method of carrying energy deep into the cloud. At larger optical depths, the effect is not as strong due to the fact that the less-dense interclump medium is capable of meaningful attenuation by itself.
The analysis above demonstrates that the angle-averaged optical depth does not provide sufficient information for describing a star-forming region. This is seen in the fact that the correlation of temperature with FOS is generally stronger than the correlation with $`\tau `$, as well as the fact that the region over which FOS is important is generally larger than the region over which $`\tau `$ is important. As a result, this suggests that indeed the ability of radiation to penetrate further into a clumpy source than a homogeneous source is due to the streaming of radiation through holes.
### 4.4 Filling factor and clump radius
In this subsection we investigate the effects of filling factor and clump radius on the dust temperature distribution. Since $`f`$ and $`r_c`$ both geometrically parameterizatize the clump distribution, we account for the range of realizations in the actual clump distribution by averaging $`T(r)`$ over nine different realizations. Finally, we note that in this subsection we present the temperature contrast ($`T(r_{\mathrm{in}})/T_{\mathrm{out}}`$), which minimizes the effect of absolute density scaling between models.
#### 4.4.1 Filling factor
Figure 12 shows a representative variation in the spherical average temperature distribution with filling factor ($`f`$) for $`r_c=0.025`$pc and $`\overline{\tau }=100`$. The temperature distributions for other models in the grid are qualitatively similar. The ratio of inner to outer temperature for the full grid (including uncertainties due to different realizations) are shown in Tables 3 and 4.
We see from these data that the role of filling factor varies with clump radius. First, the dust temperature decreases with increasing $`f`$ as expected, since a larger $`f`$ leads to a smaller FOS, and thus less heating available to the lower density medium and lower temperatures.
Second, the temperature distribution for the smallest filling factor remains essentially unchanged as the clump size changes. This is due to the high FOS throughout, so that each point is well-coupled to the external radiation field, independent of the size of the clumps.
Third, the effect of filling factor varies inversely with clump radius. This is expected, as FOS increases with increasing $`r_c`$, and decreasing $`f`$. To see this, consider the simplifying case of viewing from the center of a sphere, in which spherical clumps of radius $`r_c`$ are randomly distributed a distance $`R_{\mathrm{out}}`$ from the center, and in which no clumps overlap. This yields $`\mathrm{FOS}=1fR_{\mathrm{out}}/4r_c`$. For $`f=0.010.3`$, this yields FOS $`=0.631`$ for $`r_c=0.1`$ pc, $`0.251`$ for $`r_c=0.05`$ pc, and $`01`$ for $`r_c=0.025`$ pc. In reality FOS is influenced by two further competing factors – clump overlap in projection (increases FOS), and clump distribution in distance (decreases FOS). The actual variation of FOS with model parameters is given in Fig. 13. The result roughly agrees with the estimate above. The lower trend of the actual results imply that clump overlap is not a significant effect.
These results are somewhat different for models with high optical depth, $`\overline{\tau }=100`$. As confirmed in Table 4, the temperature distribution at these high optical depths depends significantly on $`f`$, independent of $`r_c`$.
This independence of temperature on clump radius for large optical depths can can be easily understood by considering the optical depth of individual clumps for various models. For the ISRF of Mathis, Mezger, and Panagia (1983), roughly 1/2 of the total energy is contained in $`\lambda <25\mu `$m, 2/3 in $`\lambda <115\mu `$m, and 3/4 in $`\lambda <360\mu `$m. The optical depths of a clump at these wavelengths for $`\overline{\tau }=10`$ are, $`1.1<\tau _{\mathrm{clump}}(25\mu \mathrm{m})<67`$, $`0.1<\tau _{\mathrm{clump}}(115\mu \mathrm{m})<6`$, and $`0.01<\tau _{\mathrm{clump}}(360\mu \mathrm{m})<0.8`$. On the other hand, for $`\overline{\tau }=100`$ the optical depths are ten times higher, namely, $`11<\tau _{\mathrm{clump}}(25\mu \mathrm{m})<670`$, $`1<\tau _{\mathrm{clump}}(115\mu \mathrm{m})<60`$, and $`0.1<\tau _{\mathrm{clump}}(360\mu \mathrm{m})<8`$. This implies that the clumps for $`\overline{\tau }=10`$ are opaque to 1/2 of the heating radiation, but transparent to the remaining 1/2. Consequently, $`f`$ (and FOS) can control roughly 1/2 of the heating radiation. However, for $`\overline{\tau }=100`$, the clumps are opaque to 2/3 - 3/4 of the heating radiation, meaning that $`f`$ (and FOS) can control a greater amount of the energy that penetrates to a given depth. As a result, it is not suprising to realize that as the clumps become more opaque, filling factor and the ability of radiation to stream between the clumps becomes more important.
#### 4.4.2 Clump radius
Figure 14 shows a representative variation in the spherical average temperature distribution with clump radius ($`r_c`$) for $`f=0.1`$ and $`\overline{\tau }=100`$. Again, the temperature distributions in the other models are qualitatively similar, and the temperature contrast results are shown in Tables 3 & 4.
We see from these data that the role of clump radius varies with filling factor. In comparison to the case of filling factor, the smallest clumps yield the lowest temperatures, and the largest clumps the highest temperatures, consistent with the FOS findings in Fig. 13.
The independence of temperature with $`r_c`$ for $`f=0.01`$ is due to the large FOS. In particular, there exists a large number of open lines of sight to any point, making clump variation unimportant. Conversely, as $`f`$ increases, clump size becomes more important. This is due to the fact that high filling factor models have a commensurately greater number of small clumps. As discussed previously, many small clumps are more effective at covering the sky than fewer large clumps – see Fig. 13. Consequently, clump radius is an important factor in determining the temperature profile for high filling factors.
As before, the results differ for regions of higher optical depth, $`\overline{\tau }=100`$. In Table 4 we can see the variation in the temperature contrast. As discussed previously, the increasing opacity of the clumps further amplifies the effect of filling factor and FOS on the temperature distribution, leaving the effect of clump radius as unimportant.
#### 4.4.3 Relative regions of importance
By combining the results of the two previous subsections (e.g. Fig. 13 and Tables 3 & 4), we can infer the regions of importance for filling factor and clump radius. We first consider the case of $`\overline{\tau }=10`$. At these moderate optical depths, we see that clump radius only becomes important for $`f>0.1`$. In the case of $`f<0.1`$, clump radius essentially plays no role and is dominated by the fact that the FOS is always high for such small filling factors.
On the other hand, we see that filling factor only becomes important for $`r_c<0.05`$ pc. For $`r_c>0.05`$ pc, the temperature is dominated by the fact that the FOS is relatively insensitive to $`f`$, due to the large number of holes left by scenarios with small numbers of large clumps.
It is interesting that the inferred values of $`f`$ from observations routinely fall in the range of a few per cent (e.g. Hogerheijde, Jansen, van Dishoeck 1995; Snell et al. 1984; Mundy et al. 1986; Bergin 1996). In light of the results above, this may not be a suprise. For smaller filling factors, $`f`$ and $`r_c`$ do not play a significant role in the temperature distribution. On the other hand, for larger filling factors, the cloud is closer to homogeneous, and more sensitive to clump radius than $`f`$. As a result, the temperature distribution is most sensitive to changes in $`f`$ for $`f0.1`$. While many interpretations of filling factor are not driven by continuum observations, these results are relevant since one of the dominant thermal regulators for the gas is collisions with the dust – meaning that sensitivity to $`T_{\mathrm{dust}}`$ near $`f0.1`$ can lead to sensitivity in line processes near $`f0.1`$.
As a test, we have run additional models $`f=0.03`$, and $`f=0.2`$ for the model numbers 1x21. In order to quantify the deviation from the homogenous model, we have calculated a parameter, $`\varphi ^2\frac{1}{N}(T(r_i)T_{\mathrm{hom}}(r_i))^2/\sigma (r)_{\mathrm{hom}}^2`$. Here the sum is over radial positions, $`T_{hom}(r_i)`$ is the spherical average temperature for the homogenous model, and $`\sigma (r)_{\mathrm{hom}}`$ is the uncertainty in the temperature distribution in the homogenous model. In Fig. 15 we plot the ratio of the $`\varphi (f)^2`$ to the maximum value, $`\varphi _{\mathrm{max}}^2`$ as a function of filling factor. Notice that, indeed, the deviation from the homogeneous models peaks for intermediate values of $`f`$. In particular, the greatest sensitivity to $`f`$ occurs in the range $`f=0.010.1`$ – a resultin accord with the ranges of filling factors commonly reported.
When the optical depth increases to $`\overline{\tau }=100`$, we find that filling factor plays the dominant role over a much wider range. As a result, the sensitivity to filling factor near $`f0.010.1`$ for $`\overline{\tau }=10`$ is not an effect here. At high optical depth, we also find that the clump radius is not important. This is due to the fact that $`f`$ has a larger effect on FOS than $`r_c`$ does, and that as optical depth increases, the clumps become more opaque, thereby blocking a higher fraction of the impinging radiation.
### 4.5 Range of importance of FOS
From the previous discussion, FOS is necessary to describe the local radiation field and heating within a clumpy medium. Here we attempt to infer the strength of the effect of FOS on temperature. While this has been addressed in passing, here we collect and further distill that information to help draw more general conclusions.
As seen in Section 4.3.4 and Fig. 8, the greatest effect of $`\theta (\mathrm{FOS})`$ is at large radii. We find that this is generally true across our grid of models. As expected, these points tend to have the highest FOS. Interestingly, however, the effect of FOS is on the situations where the FOS is relatively small. This can be inferred in three ways.
First, we note that the largest values of $`\theta (\mathrm{FOS})`$ occur for the models that have the lowest FOS. This can be indirectly seen in that the models with the largest $`\overline{\theta (\mathrm{FOS})}/\overline{\theta (\tau )}`$ in Table 2 tend to have the larger filling factors. As seen in Figs. 11 & 13, higher $`f`$ corresponds to lower values of FOS.
We can roughly quantify this by saying that FOS is important for those regions where FOS $`0.5`$ or so. To begin to quantify this we can consider the temperature distributions and deviations in Tables 3 & 4, together with the FOS values in Fig. 13. In particular, we note that meaningful differences in spherical average temperature occur between all models at the smallest clump radius, that essentially no variation occurs with filling factor for the largest clump radius, and that meaningful variation may occur as one changes filling factor at the intermediate clump radius. From Fig. 13 we see that the largest clump radius has values of FOS $`>0.5`$. On the other hand, for the smallest clump radius models, FOS varies from $`0.86`$ to $`0.25`$ as one changes $`f`$. Finally, for the intermediate clump radius models, as one changes from $`f=0.01`$ to $`f=0.1`$ FOS varies from $`0.95`$ to $`0.5`$, with uncertain corresponding change in the temperature distribution. But, when the difference between $`f=0.01`$ and $`f=0.3`$ is considered, there is a meaningful temperature change. Taken together, this suggests that variation in FOS for FOS $`>0.5`$ may not be as important as variations in FOS when FOS $`<0.5`$ or so.
A more direct comparison suggests that FOS always plays some role, but that a change in FOS becomes increasingly important in determining the temperature for FOS $`<0.6`$. For FOS $`>0.60.8`$, the spread in temperatures increases. To see this in Figs. 16 & 17 we have plotted the fractional deviation of the temperature of a point from the spherical average as a function of the FOS at that point for all models in our grid. The points plotted cover all radii in the models. It is clear by inspection that FOS directly correlates with the temperature deviation. Interestingly, the correlation is strongest for FOS less than $`0.6`$. On the other hand, for FOS $`>0.60.8`$, there does not appear to be a direct relationship between FOS and temperature deviation. However, for the larger values of FOS the spread in temperatures is indeed greater.
Taken together, this agrees well with the idea of streaming of photons through holes in a clumpy medium. For models with a small average FOS, a position with a few extra lines of sight can receive significant additional heating. On the other hand, when the FOS for a point is high, it is well-heated by the external radiation field. However, direct shadowing of a point by a nearby clump, and attenuation by the (assumed) tenuous interclump medium can play roles. In fact, for the high optical depth models ($`\overline{\tau }=100`$), the ranges of temperture deviations are smaller than for low optical depth models, due to the greater attenuation by the interclump medium.
## 5 Shadowing
One further impact of clumps is to directly shadow points behind them, leading to lower temperatures. This effect can be seen directly in Fig. 4. In this externally heated model, points just interior to clumps are decreased in temperature relative to equivalent points not directly behind clumps. Two important questions arise: (1) over what length scale does this effect hold?, and (2) is the cause a decrease in FOS for these points?
To answer the first question, we have considered the temperature distribution as a function of distance behind clumps. The temperature at these points is consistently lower than the spherical average by $`525`$ per cent. To identify the shadowinglengthscale, we determine the distance, $`d`$, behind the clump at which the deviation is 1/2 the maximum. A fit of $`d`$ of the form $`d=\delta \times r_c`$ yields $`\delta =1.2\pm 0.4`$. As a result, we infer shadowing is important for lengthscale of $`1.2`$ clump radii behind a clump.
To answer the second question – the role of FOS – we can reconsider Fig. 6, which plots the FOS and temperature as functions of position along three axes for a representative model. Notice how well FOS correlates with temperature, especially in the shadowed regions behind clumps. This is highly suggestive that FOS is the determining factor in shadowing. More quantitatively, we can consider the fraction of sky closed (FCS) by the clump. Using the previous notation, the fraction of closed sky is $`\mathrm{FCS}=r_c^2/(4d^2)=1/4\delta ^2`$. For $`\delta =1.2`$, this yields $`\mathrm{FCS}_{\mathrm{clump}}=0.17`$, and a corresponding FOS of $`0.83`$. This is near the limit of FOS$`=0.8`$ we inferred previously above which the sky is sufficiently open for the temperature to be insensitive to FOS. As a result, we conclude that FOS (and thus actual shadowing) is the important mechanism directly behind a clump.
## 6 Conclusions
We have constructed a grid of three-dimensional continuum radiative transfer models for clumpy star-forming regions, in order to better delineate and understand the effects of clumping on radiative transfer. Based upon this work, we find that:
1. The inclusion of clumps – even for a constant total mass / average optical depth – can significantly change the temperature distribution within the cloud. These differences in temperature can be in excess of 60 per cent, and are due to the lower effective optical depth for clumpy media relative to equivalent homogeneous media (Sect. 3).
2. The centers of clumps are warmer than would be expected on energy density grounds due to radiation trapping (Sect. 4.1).
3. The temperature distribution is driven by the ability of radiation to penetrate, and is thus strongly correlated with the angle-averaged attenuation $`e^\tau `$ (Sect. 4.2).
4. While there exists an anti-correlation of temperature with density (Sect. 4.1), the correlation with fraction of open sky (FOS) is stronger (Sect. 4.2).
5. We find only a weak correlation of dust temperature with angle-averaged optical depth, $`\tau `$ (Sect. 4.3.2). On the other hand, there exists a significant correlation between dust temperature and FOS (Sect. 4.3.3). This is interpreted as the dominance of streaming of radiation between clumps over diffusion through them in determining the radiation field (Sect. 4.3.4).
6. The dependence of radiation penetration on FOS versus $`\tau `$ is robust. The stronger correlation with FOS versus $`\tau `$ extends throughout the grid of models and for different realizations of clump distribution (Sect. 4.3.4).
7. While $`\tau `$ may be an effect near the cloud edges, FOS is more important deeper into the cloud (Sect. 4.3.4).
8. At low face-averaged optical depths, $`\overline{\tau }=10`$, filling factor is more important for small clump radii than large clump radii. At large optical depths $`\overline{\tau }=100`$, filling factor is the dominant effect for all situations (Sect. 4.4.1).
9. The effects of clump size are only important for the largest filling factors ($`f=0.3`$) and lower optical depths $`\overline{\tau }=10`$. It is unimportant for lower filling factors or larger optical depths, $`\overline{\tau }=100`$ (Sect. 4.4.2).
10. For lower face-average optical depths, $`\overline{\tau }=10`$, filling factor is most important in determining the temperature distribution for $`f=0.010.1`$, in accordance with most observations. For very opaque clouds with $`\overline{\tau }=100`$, filling factor is important over a larger range (Sect. 4.4.3).
11. FOS increases as clump radius increases and as filling factor decreases (Sect. 4.4).
12. The variation of temperature with FOS is more significant in cases of small FOS (high filling factor or small clump radii), while $`\tau `$ is relatively unimportant (Sect. 4.5).
13. For FOS $`>0.60.8`$ the sky is sufficiently open that there is little dependence of temperature on FOS (Sect. 4.5).
14. Clumps can directly shadow the regions behind them. On average, this regime extends to distance $`1.2`$ times the clump radius behind the clump, where the clump only subtends a small fraction of the sky (Sect. 5).
## Acknowledgements
We are grateful to Sheila Everett, Lee Mundy, and Dan Homan for thoughtful comments and interesting discussions, and the referee whose comments significantly improved the presentation. This work was partially supported under a grant from The Research Corporation (SDD), and Battelle.
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# HU-EP-05/28 Relativistic effects in the processes of heavy quark fragmentation
## I Introduction
The decay and production processes of the bound states with heavy quarks are investigated with greater intensity in the last years. The research aims of many experiments (ALEPH, DELPHI, SLD, CLEO, Belle, SELEX, LHC-b) are directed on the growth of experimental accuracy in the derivation of the static characteristics of heavy hadrons, their production and decay rates in different reactions BFY ; RunII ; Z00 ; QWG . The production of heavy mesons and baryons via heavy quark fragmentation in the $`e^+e^{}`$ \- annihilation represents one of the possible mechanisms for the formation of heavy hadrons with two heavy quarks. The fragmentation cross sections for the production of heavy hadrons can be calculated in an analytical form using the factorization hypothesis. The heavy quark production amplitude can be calculated on the basis of perturbative QCD. The characteristic quark virtualities of heavy quarks in the hard production are of the order of their masses while the quark virtualities in the bound state are much less than their masses due to the nonrelativistic motion. So, the total amplitude can be represented as a convolution of the hard transition amplitude with a nonperturbative factor (the wave function) determining the transition of free heavy quarks into a bound state LSG . The fragmentation mechanism was used for the study of the production processes of heavy mesons and baryons in $`e^+e^{}`$ annihilation in Refs.B1 ; CC ; B2 ; B3 ; UFN1 ; UFN2 ; MS1 ; MS2 (a more complete list of references can be found in Refs.RunII ; QWG ; UFN2 ). The growth of theoretical accuracy for the calculation of corresponding production cross sections can be reached in two ways. Firstly, it is necessary to take into account radiative corrections to the perturbative amplitude describing the production of free heavy quarks via heavy quark fragmentation. Secondly, we must consistently consider the relativistic corrections in the fragmentation amplitude connected with the relative motion of heavy quarks forming heavy hadron. From the point of view of NRQCD both effects are caused by the matrix elements as a function of the typical heavy quark velocity in the bound state rest frame of orders $`O(v_Q)`$ and $`O(v_Q^2)`$, respectively B3 . The experimental data indicate that the calculations of different production probabilities for heavy quarkonium and double heavy baryons should be improved by a systematic account of relativistic corrections. Such effects as the relative motion of heavy quarks forming heavy quarkonia and diquarks, the diquark structure effects in the calculation of the fragmentation functions of heavy diquarks should be considered. The role of relativistic effects was studied already in the processes of c-quark fragmentation into $`J/\mathrm{\Psi }`$, $`\eta _c`$ in Ref.Bashir on the basis of the Bethe-Salpeter approach, the gluon fragmentation into S-wave quarkonium in Ref.Bodwin and in the inclusive production of polarized $`J/\mathrm{\Psi }`$ from $`b`$-quark decay in Ref.Ma . The consideration of the intrinsic motion of quarks forming the heavy measons can explain the discreapancy between theoretical predictions and experimental data for the cross sections of the process $`e^+e^{}\mathrm{\Psi }\eta _c`$ BC ; BLL . The aim of the present work is to get a systematically improved description of the relativistic effects in the processes of the heavy quark fragmentation in the quasipotential approach savrin . Our goal also consists in the calculation of the relativistic corrections in heavy quark $`b`$ and $`c`$ fragmentation functions into pseudoscalar and vector heavy mesons $`(Q_1\overline{Q}_2)`$ on the basis of the relativistic quark model used earlier in the calculation of mass spectra of heavy mesons and baryons and their decay rates in different reactions rqm1 ; rqm2 ; rqm3 . In particular, we investigate double distribution probabilities for the heavy quark fragmentation over longitudinal meson momentum $`z`$ and transverse meson momentum $`p_T`$. Analytical expressions for the fragmentation probabilities as functions of transverse momentum of heavy mesons $`(Q_1\overline{Q}_2)`$ are obtained.
## II General formalism
The heavy meson production through the heavy quark fragmentation is shown in Fig.1. On the first stage $`Z^0`$ boson decays into a quark-antiquark pair with four-momenta $`q`$ and $`q^{}`$ respectively. After that one heavy quark with the four-momentum $`q`$ fragments to the heavy quarkonium. In the quasipotential approach the invariant transition amplitude of a heavy quark $`b`$ or $`c`$ into a heavy meson can be expressed as a simple convolution of a perturbative production amplitude $`T(p_1^{},p_2^{},p^{},q^{})`$ of free quarks and the quasipotential wave function of the bound state $`(Q_1\overline{Q}_2)`$ $`\mathrm{\Psi }_P(𝐩)`$ savrin ; faustov1973 :
$$M(q,P,p^{},q^{})=\frac{d𝐩}{(2\pi )^3}\overline{\mathrm{\Psi }}_P(𝐩)T(p_1^{},p_2^{},p^{},q^{}),$$
(1)
where four-momenta of fragmenting quarks $`(b,c)`$ and spectator antiquarks $`(\overline{b},\overline{c})`$ forming the heavy meson are defined as follows:
$$p_1^{}=\eta _1P+p,p_2^{}=\eta _2Pp,$$
(2)
$`p^{}`$ is the four-momentum of a free spectator quark $`b`$ or $`c`$ and $`P`$ is the four-momentum of the heavy meson. The coefficients $`\eta _{1,2}`$ in the definition (2) are taken in such a way that the following orthogonality condition is fulfilled:
$$(pP)=0,\eta _{1,2}=\frac{M^2m_{2,1}^2+m_{1,2}^2}{2M^2},$$
(3)
$`M=(m_1+m_2+W)`$ is the bound state mass.
The transition of the pair of a heavy quark and antiquark into color-singlet mesons can be envisioned as a complicated process in which the colors and spins of the heavy quark and antiquark play an important role. Different color-spin nonperturbative factors entering the amplitude $`T(p_1,p_2,p^{},q^{})`$ control the production of the heavy quark bound states. In this process the gluon virtuality $`k^2\mathrm{\Lambda }_{QCD}^2`$ and the strong coupling constant $`\alpha _s(k^2)1`$. Then the hard part of the fragmentation amplitude (1) in the leading order over $`\alpha _s`$ takes the form:
$$T(p_1^{},p_2^{},p^{},q^{})=\frac{4\alpha _s}{3\sqrt{3}}\frac{D_{\lambda \sigma }(k)}{(sm_1^2)}\overline{u}_1(p_1^{})\gamma _\lambda (\widehat{q}+m_1)\mathrm{\Gamma }v_1(q^{})\overline{u}_2(p^{})\gamma _\sigma v_2(p_2^{}),$$
(4)
where $`\mathrm{\Gamma }`$ is the vertex function for the transition of the $`Z^0`$ boson into the quark-antiquark pair; the gluon propagator is taken in the axial gauge with four-vector $`n=(1,0,0,1)`$:
$$D_{\lambda \sigma }(k)=\frac{1}{k^2+iϵ}\left[g_{\lambda \sigma }+\frac{k_\sigma n_\lambda +k_\lambda n_\sigma }{kn}\right],$$
(5)
$`s=q^2`$, $`k=(q\eta _1Pp)`$=$`(\eta _2Pp+p^{})`$ is the gluon four momentum. The color factor $`(T^a)_{il}(T^a)_{mj}\delta _{ij}/\sqrt{3}`$ = $`4\delta _{ml}/3\sqrt{3}`$ was already extracted in the amplitude (4). The transformation law of the bound state wave functions from the rest frame to the moving one with four-momenta $`P`$ is given by rqm2 ; rqm3 ; faustov1973
$$\mathrm{\Psi }_P^{\rho \omega }(𝐩)=D_1^{1/2,\rho \alpha }(R_{L_P}^W)D_2^{1/2,\omega \beta }(R_{L_P}^W)\mathrm{\Psi }_0^{\alpha \beta }(𝐩),$$
(6)
$$\overline{\mathrm{\Psi }}_P^{\lambda \sigma }(𝐩)=\overline{\mathrm{\Psi }}_0^{\epsilon \tau }(𝐩)D_1^{+1/2,\epsilon \lambda }(R_{L_P}^W)D_2^{+1/2,\tau \sigma }(R_{L_P}^W),$$
where $`R^W`$ is the Wigner rotation, $`L_P`$ is the Lorentz boost from the meson rest frame to a moving one, and the rotation matrix $`D^{1/2}(R)`$ is defined by
$$\left(\genfrac{}{}{0pt}{}{\mathrm{1\; \hspace{0.17em}0}}{\mathrm{0\; \hspace{0.17em}1}}\right)D_{1,2}^{1/2}(R_{L_P}^W)=S^1(𝐩_{1,2})S(𝐏)S(𝐩),$$
(7)
where
$$S(𝐩)=\sqrt{\frac{ϵ(p)+m}{2m}}\left(1+\frac{(𝜶𝐩)}{ϵ(p)+m}\right)$$
is the usual Lorentz transformation matrix of the four-spinor. For further transformations of the amplitude (4) the following relations are useful:
$$S_{\alpha \beta }(\mathrm{\Lambda })u_\beta ^\lambda (p)=\underset{\sigma =\pm 1/2}{}u_\alpha ^\sigma (\mathrm{\Lambda }p)D_{\sigma \lambda }^{1/2}(R_{\mathrm{\Lambda }p}^W),$$
(8)
$$\overline{u}_\beta ^\lambda (p)S_{\beta \alpha }^1(\mathrm{\Lambda })=\underset{\sigma =\pm 1/2}{}D_{\lambda \sigma }^{+1/2}(R_{\mathrm{\Lambda }p}^W)\overline{u}_\alpha ^\sigma (\mathrm{\Lambda }p).$$
Using also the transformation law of the Dirac bispinors to the rest frame
$`\overline{u}_1(𝐩)=\overline{u}_1(0){\displaystyle \frac{(\widehat{p}_1+m_1)}{\sqrt{2ϵ_1(𝐩)(ϵ_1(𝐩)+m_1)}}},p_1=(ϵ_1,𝐩),`$ (9)
(10)
$`v_2(𝐩)={\displaystyle \frac{(\widehat{p}_2m_2)}{\sqrt{2ϵ_2(𝐩)(ϵ_2(𝐩)+m_2)}}}v_2(0),p_2=(ϵ_2,𝐩),`$ (11)
we can introduce the projection operators $`\widehat{\mathrm{\Pi }}^{P,V}`$ onto the states $`(Q_1\overline{Q}_2)`$ in the meson with total spin 0 and 1 as follows:
$$\widehat{\mathrm{\Pi }}^{P,V}=[v_2(0)\overline{u}_1(0)]_{S=0,1}=\gamma _5(\widehat{ϵ}^{})\frac{1+\gamma ^0}{2\sqrt{2}}.$$
(12)
As a result the heavy quark $`b(c)`$ fragmentation amplitude into the mesons $`(b\overline{c})`$, $`(b\overline{b})`$ or $`(c\overline{c})`$ takes the form:
$$M(q,P,p^{},q^{})=\frac{2\alpha _s\sqrt{2M}}{3\sqrt{6}}\frac{D_{\lambda \sigma }(k)}{(sm_1^2)}\frac{d𝐩}{(2\pi )^3}\overline{\mathrm{\Psi }}_0(𝐩)\overline{u}_2(p^{})\gamma _\sigma \frac{(\widehat{\stackrel{~}{p}}_2m_2)}{\sqrt{2ϵ_2(𝐩)(ϵ(𝐩)+m_2)}}\times $$
(13)
$$\times \widehat{\stackrel{~}{ϵ}}^{}(v)(\widehat{v}+1)\frac{(\widehat{\stackrel{~}{p}}_1+m_1)}{\sqrt{2ϵ_1(𝐩)(ϵ_1(𝐩)+m_1)}}\gamma _\lambda (\widehat{q}+m_1)\mathrm{\Gamma }_\alpha v_1(q^{}),$$
where the four-vectors $`\stackrel{~}{ϵ}`$, $`\stackrel{~}{p}_{1,2}`$ are given by:
$$\stackrel{~}{ϵ}=L_P(0,\mathit{ϵ})=(\mathit{ϵ}𝐯,\mathit{ϵ}+\frac{(\mathit{ϵ}𝐯)𝐯}{1+v^0}),$$
(14)
$$\widehat{\stackrel{~}{p}}_{1,2}=S(L_P)\widehat{p}_{1,2}S^1(L_P),S(L_P)(1\pm \gamma ^0)S^1(L_P)=\pm (\widehat{v}\pm 1),\widehat{v}=\frac{\widehat{P}}{M}.$$
Transforming the bispinor contractions in the numerator of the expression (11) we can find the following expression for the heavy quark fragmentation amplitude including the effects of relative motion of the heavy quarks:
$$M(q,P,p^{},q^{})=\frac{2\alpha _s\sqrt{2M}}{3\sqrt{6}}\frac{\overline{\mathrm{\Psi }}_0(𝐩)}{\sqrt{\frac{ϵ_1(𝐩)}{m_1}\frac{(ϵ_1(𝐩)+m_1)}{2m_1}}\sqrt{\frac{ϵ_2(𝐩)}{m_2}\frac{(ϵ_2(𝐩)+m_2)}{2m_2}}}\frac{D_{\lambda \sigma }(k)}{(sm_1^2)}\frac{d𝐩}{(2\pi )^3}\times $$
(15)
$$\times \overline{u}_2(p^{})\gamma _\sigma [\frac{\widehat{v}1}{2}+\widehat{v}\frac{𝐩^2}{2m_2(ϵ_2+m_2)}\frac{\widehat{\stackrel{~}{p}}}{2m_2}]\widehat{\stackrel{~}{ϵ}}^{}(v)(\widehat{v}+1)\times $$
$$\times \left[\frac{\widehat{v}+1}{2}\widehat{v}\frac{𝐩^2}{2m_1(ϵ_1+m_1)}+\frac{\widehat{\stackrel{~}{p}}}{2m_1}\right]\gamma _\lambda (\widehat{q}+m_1)\mathrm{\Gamma }_\alpha v_1(q^{}).$$
The fragmentation amplitude (13) keeps at least two sources of relativistic corrections. The corrections of the first group appear from the quark-antiquark interaction operator. They can be taken into account by means of the numerical solution of the Schroedinger equation with the relevant potential. The second part of these corrections is determined by several functions depending on the momenta of the relative motion of quarks $`𝐩`$. In the limit of zero relative momentum $`𝐩`$ the amplitude $`M(q,P,p^{},q^{})`$ was studied in Ref.B1 ; CC ; B2 . At last, there exist the one-loop corrections to the fragmentation amplitude which we have not considered here.
Heavy quarkonium can be characterized by the hard momentum scale $`m`$ (the mass m of the heavy quarks), the soft momentum scale $`mv_Q`$ and the ultrasoft momentum scale $`mv_Q^2`$. We assume that the heavy quarkonium is a nonrelativistic system. This implies that the ratio $`𝐩^2/m^2v_Q^21`$. So, we introduce the expansion of all factors in Eq.(13) over relative momentum $`𝐩`$ up to terms of the second order:
$$\frac{1}{k^2}=\frac{1}{k_0^2}+\frac{1}{k_0^4}[2qpp^2]+\frac{4}{k_0^6}(qp)^2,k_0^2=\eta _2s+\eta _1m_2^2\eta _1\eta _2M^2,$$
(16)
$$\frac{1}{kn}=\frac{1}{(qn\eta _1Pn)}+\frac{pn}{(qn\eta _1Pn)^2}+\frac{(pn)^2}{(qn\eta _1Pn)^3},$$
$$(\widehat{\stackrel{~}{p}}_1+m_1)=m_1(\widehat{v}+1)+\widehat{v}\frac{𝐩^2}{2m_1}+\widehat{\stackrel{~}{p}},(\widehat{\stackrel{~}{p}}_2m_2)=m_2(\widehat{v}1)+\widehat{v}\frac{𝐩^2}{2m_2}\widehat{\stackrel{~}{p}},$$
$$\stackrel{~}{p}=L_P(0,𝐩)=(\mathrm{𝐩𝐯},𝐩+\frac{𝐯(\mathrm{𝐩𝐯})}{v^0+1}).$$
Substituting the expansions (14) into Eq.(13) we obtain:
$$M(q,P,p^{},q^{})=\frac{2\alpha _s\sqrt{2M}}{3\sqrt{6}}\overline{\mathrm{\Psi }}_0(𝐩)[1\frac{3}{8}\frac{𝐩^2}{8m_1^2}\frac{3}{8}\frac{𝐩^2}{8m_2^2}]\frac{1}{(sm_1^2)}\frac{d𝐩}{(2\pi )^3}\times $$
(17)
$$\times \overline{u}_2(p^{})\gamma _\sigma [\frac{\widehat{v}1}{2}+\widehat{v}\frac{𝐩^2}{4m_2^2}\frac{\widehat{\stackrel{~}{p}}}{2m_2}]\widehat{\stackrel{~}{ϵ}}^{}(v)(\widehat{v}+1)\times $$
$$\times [\frac{\widehat{v}+1}{2}+\widehat{v}\frac{𝐩^2}{4m_1^2}+\frac{\widehat{\stackrel{~}{p}}}{2m_1}]\gamma _\lambda (\widehat{q}+m_1)\mathrm{\Gamma }_\alpha v_1(q^{})\times $$
$$\times [\frac{1}{k_0^2}+\frac{1}{k_0^4}[2qpp^2]+\frac{4}{k_0^6}(qp)^2]\{g_{\lambda \sigma }+(k_\sigma n_\lambda +k_\lambda n_\sigma )\times $$
$$\times [\frac{1}{(qn\eta _1Pn)}+\frac{pn}{(qn\eta _1Pn)^2}+\frac{(pn)^2}{(qn\eta _1Pn)^3}]\}.$$
Let us emphasize that for the system of two heavy quarks relative motion corrections entering in the gluon propagator or heavy quark propagators have the same order $`O(v_Q^2)`$ contrary to the system including heavy and light quarks. In the last case leading order relativistic corrections are determined by the relativistic factors belonging to the light quark but other terms contain an additional small ratio $`m_q/m_Q`$. The obtained relation (15) which has the form of a three dimensional integral in the momentum space is valid when the integration is restricted to the soft momentum region, where the wave function has a significant support. Otherwise it would diverge at high momenta. Moreover, our aim consists in preserving here only the terms of the second order over $`|𝐩|/m`$ omitting corrections of higher order.
## III Heavy quark fragmentation functions into P- and V-mesons
We use the NRQCD factorization approach to the calculation of the fragmentation reactions which was developed in Refs.B1 ; B2 . The fragmentation function of the heavy quark $`Q_1`$ to produce $`{}_{}{}^{1}S_{0}^{}`$ or $`{}_{}{}^{3}S_{1}^{}`$ $`(Q_1\overline{Q}_2)`$ meson states is determined by the following expression:
$$D_{Q_1V(Q_1\overline{Q}_2)}(z)=\frac{1}{16\pi ^2}𝑑s\theta \left(s\frac{M^2}{z}\frac{m_2^2}{1z}\right)\underset{q_0\mathrm{}}{lim}\frac{|M|^2}{|M_0|^2},$$
(18)
where $`q_0`$ is the energy of the fragmentating quark: $`q_\mu =(q_0,0,0,\sqrt{q_0^2s})`$; $`M_0=\overline{u}_1(q)\mathrm{\Gamma }v_1(q^{})`$ is the amplitude of free quark $`Q_1`$ production on the mass shell, $`z`$ is the meson longitudinal momentum fraction relative to the fragmenting heavy quark. Let us consider the fragmentation production of the vector meson. Omitting the momentum of the relative motion of heavy quarks $`𝐩`$ in Eq.(14) we obtain the fragmentation amplitude which contains the leading order contribution and the correction due to the quark bound state energy $`W`$ ($`M=m_1+m_2+W`$):
$$M_1=\frac{2\sqrt{2M}\alpha _s\overline{\mathrm{\Psi }}(0)}{3\sqrt{6}}\frac{1}{(sm_1^2)(\eta _2s+\eta _1m_2^2\eta _1\eta _2M^2)}\times $$
(19)
$$\times [\overline{u}_2(p^{})2\widehat{\stackrel{~}{ϵ}}^{}(v)(\widehat{q}+m_1)\mathrm{\Gamma }_\alpha v_1(q^{})+\frac{(s+\eta _2m_1M\eta _1M^2)}{(nq\eta _1nP)}\overline{u}_2(p^{})\widehat{n}\widehat{\stackrel{~}{ϵ}}^{}(v)(\widehat{v}+1)\mathrm{\Gamma }_\alpha v_1(q^{})+$$
$$+\frac{(m_1\eta _1M)}{(nq\eta _1nP)}\overline{u}_2(p^{})\widehat{n}\widehat{\stackrel{~}{ϵ}}^{}(v)(\widehat{v}+1)\widehat{p}^{}\mathrm{\Gamma }_\alpha v_1(q^{})+\frac{(m_2\eta _2M)}{(nq\eta _1nP)}\overline{u}_2(p^{})\widehat{n}\widehat{\stackrel{~}{ϵ}}^{}(v)(\widehat{v}+1)(\widehat{q}+m_1)\mathrm{\Gamma }_\alpha v_1(q^{})].$$
Taking into account linear terms in the binding energy $`W`$ in the expansion of Eq.(16) we next perform an averaging and summation over the meson spins in initial and final states in the square modulus $`|M|^2`$
$$\frac{1}{3}\underset{spin}{}\stackrel{~}{ϵ}_\alpha ^{}(v)\stackrel{~}{ϵ}_\beta (v)=\frac{1}{3}\left(g_{\alpha \beta }+v_\alpha v_\beta \right),$$
(20)
and then consider the limit $`q_0\mathrm{}`$ in the obtained expression. In this limit $`\widehat{P}`$ and $`\widehat{q}`$ have the order of the $`M_Z`$ mass and the coefficients in corresponding expressions are determined by the heavy quark masses of the fragmenting quark $`m_1`$ and the spectator quark $`m_2`$. In the leading order we can substitute $`P=zq`$, and the trace in $`|M|^2`$ over the Dirac indices is proportional to $`Tr(\mathrm{\Gamma }_\alpha \widehat{q^{}}\mathrm{\Gamma }_\beta \widehat{q^{}})`$. It will disappear in the ratio $`|M|^2/|M_0|^2`$. Then the fragmentation probability (16) can be written as a sum of two terms:
$$D_{Q_1V(Q_1\overline{Q}_2)}(z)=\frac{8\alpha _s^2|\mathrm{\Psi }(0)|^2}{27m_2^3}\frac{rz(1z)^2}{[1(1r)z]^6}(v_0+v_1),$$
(21)
$$v_0=22(32r)z+3(32r+4r^2)z^22(1r)(4r+2r^2)z^3+(1r)^2(32r+2r^2)z^4,$$
(22)
$$v_1=\frac{W}{3m_2[1(1r)z]^2}[12+6r+(6084r+48r^2)z+(138+259r296r^2+102r^3)z^2+$$
(23)
$$+(192420r+552r^2536r^3+330r^4)z^3+(168+446r614r^2+746r^3612r^4+232r^5)z^4+$$
$$+(84274r+476r^2650r^3+574r^4272r^5+62r^6)z^5+(18+57r126r^2+282r^3$$
$$412r^4+333r^5144r^6+28r^7)z^6+r(1040r+56r^224r^314r^4+16r^54r^6)z^7].$$
As mentioned above there exist several sources of relativistic corrections $`|𝐩|^2/m_{1,2}^2`$ in the expression (15). The first part of terms appears from the expansion of the gluon propagator and relativistic factors in the Dirac bispinors. The structure of the spinor matrix element in this case is the same as in Eq.(17):
$$M_{21}(q,P,p^{},q^{})=\frac{2\alpha _s\sqrt{2M}}{3\sqrt{6}}\overline{\mathrm{\Psi }}_{(\overline{Q}_1Q_2),0}(𝐩)\frac{d𝐩}{(2\pi )^3}\frac{(m_1+m_2)}{m_2(sm_1^2)^2}\times $$
(24)
$$\times \{[1\frac{p^2}{k_0^2}+\frac{4(qp)^2}{k_0^4}\frac{𝐩^2}{8m_1^2}\frac{𝐩^2}{8m_2^2}]\overline{u}_2(p^{})2\widehat{\stackrel{~}{ϵ}}^{}(v)(\widehat{q}+m_1)\mathrm{\Gamma }_\alpha v_1(q^{})+$$
$$+[1\frac{p^2}{k_0^2}+\frac{4(qp)^2}{k_0^4}\frac{𝐩^2}{8m_1^2}\frac{𝐩^2}{8m_2^2}+\frac{2(qp)(pn)}{k_0^2(nq\eta _1nP)}+\frac{(pn)^2}{(nq\eta _1nP)^2}]\times $$
$$\times \frac{(sm_1^2)}{(nq\eta _1nP)}\overline{u}_2(p^{})\widehat{n}\widehat{\stackrel{~}{ϵ}}^{}(v)(\widehat{v}+1)\mathrm{\Gamma }_\alpha v_1(q^{})\}.$$
Another part of corrections is determined both by the gluon propagator terms and relativistic addenda $`\widehat{\stackrel{~}{p}}`$ in the square brackets of Eq.(15):
$$M_{22}(q,P,p^{},q^{})=\frac{2\alpha _s\sqrt{2M}}{3\sqrt{6}}\overline{\mathrm{\Psi }}_{(\overline{Q}_1Q_2),0}(𝐩)\frac{d𝐩}{(2\pi )^3}\frac{(m_1+m_2)}{m_2(sm_1^2)^2}\times $$
(25)
$$\times \left[\frac{(np)}{(nq\eta _1nP)}+\frac{2qp}{k_0^2}\right]\frac{(p_\sigma n_\lambda +p_\lambda n_\sigma )}{(nq\eta _1nP)}\overline{u}_2(p^{})\gamma _\sigma \widehat{\stackrel{~}{ϵ}}^{}(v)(\widehat{v}+1)\gamma _\lambda (\widehat{q}+m_1)\mathrm{\Gamma }_\alpha v_1(q^{}),$$
$$M_{23}(q,P,p^{},q^{})=\frac{2\alpha _s\sqrt{2M}}{3\sqrt{6}}\overline{\mathrm{\Psi }}_{(\overline{Q}_1Q_2),0}(𝐩)\frac{d𝐩}{(2\pi )^3}\frac{(m_1+m_2)}{m_2(sm_1^2)^2}\times $$
(26)
$$\times \frac{(p_\sigma n_\lambda +p_\lambda n_\sigma )}{(nq\eta _1nP)}\overline{u}_2(p^{})\gamma _\sigma \left[\frac{\widehat{\stackrel{~}{p}}}{2m_2}\widehat{\stackrel{~}{ϵ}}^{}(v)(\widehat{v}+1)+\widehat{\stackrel{~}{ϵ}}^{}(v)(\widehat{v}+1)\frac{\widehat{\stackrel{~}{p}}}{2m_1}\right]\gamma _\lambda (\widehat{q}+m_1)\mathrm{\Gamma }_\alpha v_1(q^{}),$$
$$M_{24}(q,P,p^{},q^{})=\frac{2\alpha _s\sqrt{2M}}{3\sqrt{6}}\overline{\mathrm{\Psi }}_{(\overline{Q}_1Q_2),0}(𝐩)\frac{d𝐩}{(2\pi )^3}\frac{(m_1+m_2)}{m_2(sm_1^2)^2}\times $$
(27)
$$\times \left[g_{\lambda \sigma }\frac{(k_\sigma n_\lambda +k_\lambda n_\sigma )}{(nq\eta _1nP)}\right]\overline{u}_2(p^{})\gamma _\sigma \frac{\widehat{\stackrel{~}{p}}}{2m_2}\widehat{\stackrel{~}{ϵ}}^{}(v)(\widehat{v}+1)\frac{\widehat{\stackrel{~}{p}}}{2m_1}\gamma _\lambda (\widehat{q}+m_1)\mathrm{\Gamma }_\alpha v_1(q^{}).$$
Further transformations of the interference terms $`(M_{21}M_{2i}^{}+M_{2i}M_{21}^{})`$ ($`i=2,3,4`$) in the square modulus $`|M|^2`$ are performed by means of the system Reduce reduce . The relation
$$𝑑\mathrm{\Omega }_𝐩p_\mu p_\nu =(g_{\mu \nu }v_\mu v_\nu )\frac{𝐏^23P^{02}}{9P^{02}}𝑑\mathrm{\Omega }_𝐩𝐩^2,$$
(28)
which accounts the orthogonality condition (3) is used for the integration over the angle variables in the relative momentum space. Omitting numerous intermediate expressions appearing in the calculation $`|M|^2`$ which have sufficiently cumbersome forms we present the final result for the relativistic correction in the fragmentation function (16):
$$D_{Q_1V(Q_1\overline{Q}_2)}^{rel}(z)=\frac{8\alpha _s^2|\mathrm{\Psi }(0)|^2}{27m_2^3}\frac{rz(1z)^2}{[1(1r)z]^6}v_2,$$
(29)
$$v_2(z)=\frac{𝐩^2}{36[1(1r)z]^2(1r)^2m_2^2}[2+16r(23r)+2(12r)(544r+30r^2)z+$$
(30)
$$+(23+18r206r^2+508r^3432r^4)z^2+8(1r)(4+40r11r^250r^3+66r^4)z^3$$
$$2(1r)^2(14+214r+47r^2248r^3+120r^4)z^4+2(1r)^3(7+99r+120r^2230r^3+108r^4)z^5$$
$$(1r)^4(3+26r+82r^284r^3+48r^4)z^6],$$
here $`𝐩^2`$ is the special parameter determining the numerical value of relativistic effects (see the discussion before Eq.(60)). Integrating expressions (19), (27) over $`z`$ we obtain the total fragmentation probability:
$$\mathrm{\Omega }_V=_0^1D_{Q_1V(Q_1\overline{Q}_2)}(z)𝑑z=\frac{8\alpha _s^2|\mathrm{\Psi }(0)|^2}{405m_2^3(1r)^6}\left(f_{0v}(r)+\frac{W}{m_2}f_{1v}(r)+\frac{𝐩^2}{m_2^2}f_{2v}(r)\right),$$
(31)
$$f_{0v}(r)=24+85r235r^2+300r^385r^489r^5+15r(74r+3r^2+10r^3+2r^4)\mathrm{ln}r,$$
(32)
$$f_{1v}(r)=\frac{1}{42(1r)^2}[2607+7185r+23576r^2116018r^3+159670r^4170373r^5+153860r^6$$
(33)
$$34906r^720387r^8+210r(42+223r388r^2+236r^3268r^443r^5+212r^6+28r^7)\mathrm{ln}r],$$
$$f_{2v}(r)=\frac{1}{252(1r)^2}[4883017r+48979r^2201740r^3+136955r^4+23597r^520958r^6+15696r^7$$
(34)
$$105r(7200r+347r^2+1194r^3222r^4+156r^5+48r^6)\mathrm{ln}r].$$
The parameter $`r=m_2/(m_1+m_2)`$ is the ratio of the constituent mass of a spectator quark to the mass of two quarks composing the heavy meson. The contributions to the fragmentation functions (20)-(21) and (28) give the distributions in the longitudinal momentum $`z`$ of the heavy mesons. They are shown in Fig.2. We find it convenient to present here the fragmentation functions of antiquarks for the comparison with experimental data. To extend the present calculations to the distributions in the transverse momentum $`p_T`$ of the the heavy meson, the following relation between the invariant mass $`s`$ of the fragmenting heavy quark, the transverse momentum $`p_T`$ and longitudinal momentum $`z`$ is taken into account Close :
$$s(z,t)=\frac{M^2+p_T^2}{z}+\frac{m_2^2+p_T^2}{1z}.$$
(35)
Introducing further the dimensionless variable $`t=p_T/(m_1+m_2)`$, we can determine the $`p_T`$-dependent relativistic corrections to the fragmentation probability $`D_{Q_1V(Q_1\overline{Q}_2)}(t)`$ B4 :
$$D_{Q_1V(Q_1\overline{Q}_2)}(t)=_0^1𝑑z\frac{2M^2t}{z(1z)}D_{Q_1V(Q_1\overline{Q}_2)}(z,s(z,t))=$$
(36)
$$=\frac{4\alpha _s^2|\mathrm{\Psi }(0)|^2r}{27m_2^3(1r)^6t^6}\left[D_{0v}(t)+\frac{W}{m_2}D_{1v}(t)+\frac{𝐩^2}{m_2^2}D_{2v}(t)\right],$$
$$D_{0v}(t)=\{(30r^3+30r^4)t+(61r33r^248r^3+20r^4)t^3+(5+13r16r^2+4r^3+4r^4)t^5+$$
(37)
$$+\mathrm{arctan}\frac{t(1r)}{r+t^2}[30r^4+(99r^266r^3+30r^4)t^2+(9+20r+r^2+22r^3+8r^4)t^4+$$
$$+(912r+4r^2+8r^3)t^6]+\mathrm{ln}(r)[96r^3t+(48r+56r^2+16r^3)t^3]+$$
$$+\mathrm{ln}\frac{(1+t^2)}{(r^2+t^2)}[48r^3t+(24r+28r^2+8r^3)t^3+(4+8r4r^2)t^7]\},$$
$$D_{1v}(t)=\frac{1}{3t^2(1r)^2}\{\frac{24rt(1r)^6}{1+t^2}+(24r144r^2+360r^3480r^4+45r^5+486r^6291r^7)t+$$
(38)
$$+(24r+144r^2633r^3+2115r^41233r^51266r^6+897r^7)t^3+$$
$$+(18r+1130r^24129r^3+3938r^4+595r^51994r^6+478r^7)t^5+$$
$$+(30215r+581r^2742r^3+738r^4596r^5+148r^6+56r^7)t^7+$$
$$+(60r255r^2+360r^3195r^4+30r^5)t^9]+$$
$$+\mathrm{arctan}\frac{t(1r)}{(r+t^2)}[315r^6315r^7+(90r^41650r^5345r^6+960r^7)t^2+$$
$$+(63r^21800r^3+5646r^4714r^53495r^6+741r^7)t^4+$$
$$+(54+237r33r^2807r^3+642r^4621r^5+423r^6+150r^7)t^6+$$
$$+(54+297r507r^2+216r^3252r^4+147r^5+126r^6)t^8+$$
$$+(60r+195r^2165r^3+30r^4)t^{10}]+$$
$$+\mathrm{ln}(r)[(576r^5+288r^7)t+(432r^3+48r^4+5856r^54416r^648r^7)t^3+$$
$$+(144r+1296r^22136r^3768r^4+1320r^5+432r^6)t^5]+$$
$$+\mathrm{ln}\frac{(1+t^2)}{(r^2+t^2)}[(288r^5+144r^7)t+(216r^3+24r^4+2928r^52208r^624r^7)t^3+$$
$$+(72r+648r^21068r^3384r^4+660r^5+216r^6)t^5+$$
$$+(24156r+312r^2192r^3+72r^460r^5)t^9]\},$$
$$D_{2v}(t)=\frac{1}{36t^2(1r)^2}\{420r^5(1r)t+(6074r^320646r^4+11976r^5+2596r^6)t^3+$$
(39)
$$+(573r+2309r^2+5008r^35180r^41636r^5+72r^6)t^5+$$
$$+(5189r+1124r^2960r^3+68r^4+48r^596r^6)t^7+$$
$$+\mathrm{arctan}\frac{t(1r)}{r+t^2}[420r^6+(7470r^4+15480r^5+5040r^6)t^2+$$
$$+(1887r^21758r^320334r^46936r^5+900r^6)t^4+$$
$$+(9+100r+1303r^21614r^3632r^4+216r^5168r^6)t^6+$$
$$+(9+108r1336r^296r^3+64r^4192r^5)t^8]+$$
$$+\mathrm{ln}(r)[(3072r^5+4608r^6)t+(6240r^310688r^424256r^5+768r^6)t^3+$$
$$+(240r+712r^2+5232r^3+3392r^4+384r^5)t^5]+$$
$$+\mathrm{ln}\frac{(r^2+t^2)}{(t^2+1)}[(1536r^5+2304r^6)t+(3120r^35344r^412128r^5+384r^6)t^3+$$
$$+(120r+356r^2+2616r^3+1696r^4+192r^5)t^5+$$
$$+(4+72r+324r^2208r^3+96r^4)t^9]\}.$$
The leading order contribution (35) coincides with the result of Ref.B4 and the two other terms (36) and (37) determine the relativistic and bound state corrections. The functions (35)-(37) are plotted in Fig.4.
Let us next consider the calculation of the fragmentation functions into pseudoscalar heavy mesons $`\eta _c`$, $`B_c`$, $`\eta _b`$. The general expression of the fragmentation amplitude consists of three terms:
$$M_3=\frac{2\sqrt{2M}\alpha _s|\mathrm{\Psi }(0)|}{3\sqrt{6}}\frac{(m_1+m_2)}{m_2(sm_1^2)^2}[a_1\overline{u}_2(p^{})\gamma _5(\widehat{q}+m_1)\mathrm{\Gamma }_\alpha v_1(q^{})+$$
(40)
$$+a_2\overline{u}_2(p^{})\gamma _5\mathrm{\Gamma }_\alpha v_1(q^{})+a_3\overline{u}_2(p^{})\gamma _5\widehat{n}(\widehat{q}+m_1)\mathrm{\Gamma }_\alpha v_1(q^{})],$$
(the fourth possible term $`\overline{u}_2(p^{})\gamma _5\widehat{n}\mathrm{\Gamma }_\alpha v_1(q^{})`$ does not contribute to the fragmentation function) where the coefficients $`a_i`$ (i=1,2,3) contain the leading order contribution and corrections proportional to $`𝐩^2`$ and $`W`$:
$$a_1=1\frac{2}{9}\frac{𝐩^2}{(sm_1^2)r}+\frac{𝐩^2}{m_1^2}[\frac{5}{24}+\frac{2}{9}r\frac{5}{9}r^2+$$
(41)
$$+\frac{1}{1z(1r)}(\frac{1}{2}zr^3\frac{2}{3}zr^2\frac{1}{6}r^2\frac{1}{18}rz+\frac{7}{18}r+\frac{2}{9}z\frac{2}{9})]+$$
$$+\frac{2(m_1+m_2)W}{(sm_1^2)}(1r)+\frac{W}{(m_1+m_2)}\left(2\frac{1}{r}\frac{z}{1z(1r)}+\frac{2rz}{1z(1r)}\right),$$
$$a_2=1+\frac{𝐩^2}{m_1^2}\left(\frac{1}{6}r^2+\frac{5}{18}r\frac{17}{12}\right)+\frac{2}{9}\frac{𝐩^2}{m_2^2}\frac{r^2z^2}{[1z(1r)]^2}+$$
(42)
$$+\frac{𝐩^2}{m_1^2}\frac{1}{1z(1r)}\left(\frac{2}{9}r^2z+\frac{2}{9}r^2+\frac{1}{3}rz\frac{1}{3}r\frac{1}{9}z+\frac{1}{9}\right)+$$
$$+\frac{𝐩^2}{m_1^2}\frac{m_2^2}{(sm_1^2)}\left(\frac{5}{9}r\frac{5}{9rz}+\frac{1}{r}+\frac{2}{9r^2z}\frac{2}{9r^2}+\frac{1}{3z}\frac{4}{3}\right)\frac{4}{9r}\frac{𝐩^2}{(sm_1^2)}+$$
$$+\frac{2(m_1+m_2)W}{(sm_1^2)}(1r)+\frac{W}{(m_1+m_2)}\left(2\frac{1}{r}\right),$$
$$a_3=1\frac{𝐩^2}{9m_1^2}\left(r+\frac{1}{8}\right)+\frac{𝐩^2}{m_1m_2}\frac{rz}{1z(1r)}+\frac{2r^2z^2}{9}\frac{𝐩^2}{m_2^2}$$
(43)
$$\frac{2}{9r}\frac{𝐩^2}{(sm_1^2)}+\frac{2(m_1+m_2)W}{(sm_1^2)}(1r)+\frac{W}{(m_1+m_2)}\left(2\frac{1}{r}\right).$$
The exact expressions (39)-(41) are obtained on the basis of Eq.(15) ($`\stackrel{~}{ϵ}\gamma _5`$) keeping the terms of the second order over relative momentum $`p`$ and binding energy corrections. The squared modulus amplitude $`|M_3|^2`$ leads by using the definition (16) after the integration over $`s`$ to the following fragmentation distribution for the production of the pseudoscalar heavy mesons:
$$D_{Q_1P(Q_1\overline{Q}_2)}(z)=\frac{8\alpha _s^2|\mathrm{\Psi }(0)|^2}{81m_2^3}\frac{rz(1z)^2}{[1(1r)z]^6}(p_0+p_1+p_2),$$
(44)
$$p_0=6+18(2r1)z+(2174r+68r^2)z^22(1r)(619r+18r^2)z^3+3(1r)^2(12r+2r^2)z^4,$$
(45)
$$p_1=\frac{W}{m_2[1(1r)z]^2}\{12+6r+60z+24r(5+2r)z+[126+r(425482r+218r^2)]z^2+$$
(46)
$$+2[72+r(329+2r(298+r(283+131r)))]z^32(1r)[48+r(219+r(410+7r(61+33r))]z^4$$
$$+2(1r)^2[18+r(79+2r(69+13r(5+3r)))]z^5+(r1)^3[6+r(25+r(35+6r(4+3r)))]z^6\},$$
$$p_2=\frac{𝐩^2}{36(1r)^2m_2^2[1(1r)z]^2}\{648r6[5+2r(3+r(3+22r))]z$$
(47)
$$[63+2r(191+r(431+4r(197+146r)))]z^2+8(1r)[9+r(108+r(217+5r(42+29r)))]z^3$$
$$2(1r)^2[24+r(390+r(697+4r(53+5r)))]z^4+2(r1)^3[9+r(155+2r\times $$
$$(122+r(7+36r)))]z^5+(1r)^4(3+42r58r^2+24r^4)z^6\}.$$
Integrating expressions (43)-(45) over $`z`$ we obtain the total fragmentation probability for the pseudoscalar mesons:
$$\mathrm{\Omega }_P=_0^1D_{Q_1P(Q_1\overline{Q}_2)}(z)𝑑z=\frac{8\alpha _s^2|\mathrm{\Psi }(0)|^2}{405m_2^3(1r)^6}\left(f_{0p}(r)+\frac{W}{m_2}f_{1p}(r)+\frac{𝐩^2}{m_2^2}f_{2p}(r)\right),$$
(48)
$$f_{0p}(r)=8+5r+215r^2440r^3+265r^453r^5+15r(1+8r+r^26r^3+2r^4)\mathrm{ln}r,$$
(49)
$$f_{1p}(r)=\frac{1}{42(1r)}[8973640r+8400r^259850r^3+147105r^4132762r^5+46830r^6$$
(50)
$$6980r^7+210r(63r129r^2+115r^3+123r^484r^5+18r^6)\mathrm{ln}r],$$
$$f_{2p}(r)=\frac{1}{756(1r)^2}[6149051r9681r^2+106470r^373815r^445129r^5+$$
(51)
$$+41146r^610554r^7+105r(3204r+281r^2+798r^3402r^4+48r^5+24r^6)\mathrm{ln}r].$$
The transverse momentum fragmentation functions for the production of pseudoscalar mesons $`\eta _c`$, $`B_c`$ and $`\eta _b`$ can be derived in a similar way as for the vector mesons. The corresponding expressions are given as follows:
$$D_{Q_1P(Q_1\overline{Q}_2)}(t)=_0^1𝑑z\frac{2M^2t}{z(1z)}D_{Q_1P(Q_1\overline{Q}_2)}(z,s(z,t))=$$
(52)
$$=\frac{4\alpha _s^2|\mathrm{\Psi }(0)|^2r}{81m_2^3(1r)^6t^6}\left[D_{0p}(t)+\frac{W}{m_2}D_{1p}(t)+\frac{𝐩^2}{m_2^2}D_{2p}(t)\right],$$
$$D_{0p}(t)=\{(1r)t[30r^3r(6120r+28r^2)t^2(348r+48r^212r^3)t^4]+$$
(53)
$$+3\mathrm{arctan}\frac{t(1r)}{r+t^2}[10r^43r^2(11+2r+2r^2)t^2+(3+4r+19r^26r^3)t^4+$$
$$+(3+12r20r^2+8r^3)t^6)]24rt\mathrm{ln}(r)[4r^2(2+r+2r^2)t^2]$$
$$12t\mathrm{ln}\frac{(1+t^2)}{(r^2+t^2)}[4r^3r(2+r+2r^2)t^2+(1r)^2t^6]\},$$
$$D_{1p}(t)=\frac{1}{t^2(1r)}\{\frac{24rt(1r)^4(2r1)}{1+t^2}+3r(848r+112r^2128r^3+37r^4+19r^5)t+$$
(54)
$$+r(24+144r691r^2+1032r^3+84r^4545r^5)t^3+$$
$$+r(10+360r1180r^2+1037r^3495r^4+268r^5)t^5$$
$$3(219r+117r^2300r^3+320r^4140r^5+20r^6)t^7+$$
$$+3\mathrm{arctan}\frac{t(1r)}{(r+t^2)}[35r^610r^4(13+15r+25r^2)t^2+$$
$$+r^2(15281r+501r^2+96r^3+79r^4)t^4+$$
$$+(6+31r7r^2+79r^3245r^4+42r^5+8r^6)t^6$$
$$(611r68r^2+208r^3160r^4+40r^5)t^8+$$
$$+24rt\mathrm{ln}(r)[16r^5+2r^2(710r+38r^2+4r^3)t^2+(2+20r21r^25r^322r^4)t^4]+$$
$$+12t\mathrm{ln}\frac{(r^2+t^2)}{1+t^2}[16r^62r^3(710r+38r^2+4r^3)t^2+$$
$$+r(220r+21r^2+5r^3+22r^4)t^4(1r)^2(25r+5r^2)t^8]\},$$
$$D_{2p}(t)=\frac{1}{36t^2(1r)^2}\{420r^5(1r)t+2r^3(19577323r+4368r^2+998r^3)t^3$$
(55)
$$r(4292629r+152r^2396r^3+1684r^4+760r^5)t^5+$$
$$+3(1+65r308r^2+360r^3108r^424r^5+16r^6)t^7+$$
$$+\mathrm{arctan}\frac{t(1r)}{r+t^2}[(420r^6+(5310r^4+11640r^5+4440r^6)t^2+$$
$$+(1023r^24158r^39486r^44080r^51764r^6)t^4+$$
$$+(9324r9r^2+1674r^3+1632r^4360r^5+72r^6)t^6+$$
$$+(936r+744r^2552r^324r^4+96r^5)t^8]+$$
$$+24rt\mathrm{ln}(r)[128r^4+192r^54r^2(35+110r+110r^2+26r^3)t^2+$$
$$+(10r+150r^2+152r^3+100r^4+4r^5)t^4]+$$
$$+12t\mathrm{ln}\frac{(1+t^2)}{(r^2+t^2)}[128r^5+192r^64r^3(35+110r+110r^2+26r^3)t^2+$$
$$+r(10r+150r^2+152r^3+100r^4+4r^5)t^4(1+14r13r^24r^3+r^4)t^8)]\}.$$
The leading order function (51) coincides with the result of Ref.B4 , and the two other functions (52) and (53) determine the relativistic and binding energy corrections. Functions (51)-(53) are plotted in Fig.5. The obtained results for the transverse momentum distributions are valid for the transverse momentum $`p_T`$ up to values of order of the meson mass. The further growth of the momentum $`p_T`$ demands the consideration of the omitted corrections of order $`O(p_T/M_Z)`$.
## IV Quasipotential quark model
To estimate numerical values of the investigated effects in the heavy quark fragmentation we used the relativistic quark model. In the quasipotential approach the bound state of a quark and antiquark is described by the Schrödinger type equation rqm4
$$\left(\frac{b^2(M)}{2\mu _R}\frac{𝐩^2}{2\mu _R}\right)\mathrm{\Psi }_0(𝐩)=\frac{d𝐪}{(2\pi )^3}V(𝐩,𝐪,M)\mathrm{\Psi }_0(𝐪),$$
(56)
where the relativistic reduced mass is
$$\mu _R=\frac{E_1E_2}{E_1+E_2}=\frac{M^4(m_1^2m_2^2)^2}{4M^3},$$
(57)
and the particle energies $`E_1`$, $`E_2`$ are given by
$$E_1=\frac{M^2m_2^2+m_1^2}{2M},E_2=\frac{M^2m_1^2+m_2^2}{2M},$$
(58)
here $`M=E_1+E_2`$ is the bound state mass, $`m_{1,2}`$ are the masses of heavy quarks ($`Q_1`$ and $`Q_2`$) which form the meson, and $`𝐩`$ is their relative momentum. In the center of mass system the relative momentum squared on mass shell reads
$$b^2(M)=\frac{[M^2(m_1+m_2)^2][M^2(m_1m_2)^2]}{4M^2}.$$
(59)
The kernel $`V(𝐩,𝐪,M)`$ in Eq. (56) is the quasipotential operator of the quark-antiquark interaction. Within an effective field theory (NRQCD) the quark-antiquark potential was constructed in Ref.NR1 ; NR2 by the perturbation theory improved by the renormalization group resummation of large logarithms. In the quasipotential quark model the kernal $`V(𝐩,𝐪,M)`$ is constructed phenomenologically with the help of the off-mass-shell scattering amplitude, projected onto the positive energy states. The heavy quark-antiquark potential with the account of retardation effects and the one loop radiative corrections can be presented in the form of a sum of spin-independent and spin-dependent parts. Explicit expressions for it are given in Refs. rqm2 ; rqm3 . Taking into account the accuracy of the calculation of relativistic corrections to the fragmentation probabilities, we can use for the description of the bound system $`(Q_1\overline{Q}_2)`$ the following simplified interaction operator in the coordinate representation:
$$\stackrel{~}{V}(r)=\frac{4}{3}\frac{\overline{\alpha }_V(\mu ^2)}{r}+Ar+B,$$
(60)
where the parameters of the linear potential $`A=0.18GeV^2`$, $`B=0.16GeV`$,
$$\overline{\alpha }_V(\mu ^2)=\alpha _s(\mu ^2)\left[1+\left(\frac{a_1}{4}+\frac{\gamma _E\beta _0}{2}\right)\frac{\alpha _s(\mu ^2)}{\pi }\right],$$
(61)
$$a_1=\frac{31}{3}\frac{10}{9}n_f,\beta _0=11\frac{2}{3}n_f.$$
Here $`n_f=3`$ is the number of flavors and $`\mu `$ is a renormalization scale. All the parameters of the model like quark masses, parameters of the linear confining potential $`A`$ and $`B`$, mixing coefficient $`\epsilon `$ and anomalous chromomagnetic quark moment $`\kappa `$ entering in the quasipotential $`V(𝐩,𝐪,M)`$ were fixed from the analysis of heavy quarkonium masses rqm1 ; rqm2 ; rqm3 and radiative decays rqm2 . The heavy quark masses $`m_b=4.88`$ GeV, $`m_c=1.55`$ GeV and the parameters of the linear potential $`A=0.18`$ GeV<sup>2</sup> and $`B=0.16`$ GeV have standard values of the quark models. Solving the Schrödinger-like quasipotential equation we obtain an initial expression for the bound state wave functions in the case of $`(c\overline{c})`$, $`(\overline{b}c)`$ and $`(\overline{b}b)`$ systems. For numerical estimations of relativistic effects in the production of heavy mesons via heavy quark fragmentation we need the values of the wave functions at the origin, the bound state energy and the parameter of relativistic effects: $`𝐩^2\overline{\mathrm{\Psi }}_0(𝐩)𝑑𝐩/(2\pi )^3`$. Note that this integral would diverge in the high momentum region. Obviously, the reason of this divergence is connected with the used expansion of the integral function in the basic equation (13) over the ratio $`𝐩^2/m^2`$. In the coordinate representation this divergence would appear when we set $`𝐫=0`$ in the Coulomb part of the potential (58). Different regularization prescriptions are commonly used in this case aKM ; Labelle ; KM ; Khan . An approach to the calculation of this integral was formulated in Ref.KM1 for solving the orthopositronium decay problem in quantum electrodynamics. Their prescription is in an agreement with the calculations carried out in the effective field theories Adkins . Unfortunately, in the investigation of the bound states of heavy quarks we cannot use it because the valid wave function asymptotics at $`p\mathrm{}`$ is not the Coulomb-like. So, to fix the value of the relativistic correction we explore the dimensional regularization scheme where the scaleless momentum integral $`V(𝐩𝐪)\mathrm{\Psi }(𝐪)\frac{d^d𝐪}{(2\pi )^d}\frac{d^d𝐩}{(2\pi )^d}`$ related to the problem vanishes aKM ; Labelle ; CMY . Then we can express the necessary quantity in the form:
$$𝐩^2\frac{1}{\mathrm{\Psi }(0)}\frac{d^d𝐩}{(2\pi )^d}𝐩^2\overline{\mathrm{\Psi }}_0(𝐩)=2\mu _R\stackrel{~}{W}+2\mu _R|B|.$$
(62)
The solutions of the Schrödinger-like equation (54) with the potential (58) determine the energy spectrum $`\stackrel{~}{W}`$ of the heavy quark system and lead to the numerical values of the parameter (60) for the bound states $`(\overline{c}c)`$, $`(\overline{b}b)`$ and $`(\overline{b}c)`$ which are presented in Table I. They are in qualitative agreement with the other possible approach for the estimation of the value (60) based on the natural regularization directly connected with the relativistic structure factors entering in the fragmentation amplitude (13). Heavy quark symmetry predicts that the wave functions of the vector and pseudoscalar states are different due to corrections of order $`v_Q^2`$. The analogous statement is valid for the parameter $`𝐩^2`$. Nevertheless, in this study we neglect this difference and write in Table I equal values for $`\mathrm{\Psi }(0)`$ and $`𝐩^2`$ for $`V`$\- and $`P`$-mesons. Our value $`𝐩^2=0.5GeV^2`$ for $`(c\overline{c})`$-states is slightly smaller than $`𝐩^2=0.7GeV^2`$ used in Ref.Bashir where it was fixed from the analysis of the quarkonium decay rates. The theoretical uncertainty of the obtained values $`𝐩^2`$ in the Table I is determined by perturbative and nonperturbative corrections to the quasipotential rqm1 ; rqm2 and is not exceeding $`30\%`$.
## V Discussion and Conclusions
As mentioned above, the problem of heavy hadron production in $`e^+e^{}`$ and $`p\overline{p}`$ collisions became very urgent in last years. The experimental investigations carried out in this field allowed to measure the $`b`$ quark fragmentation function in $`Z^0`$ decays Z00 . The study of the fragmentation processes is important as a tool to reveal the features of nonperturbative quantum chromodynamics. There appear experimental data indicating essential differences between the theoretical predictions and experiment BFY ; RunII ; QWG ; BC ; BLL . In the present study we investigated the role of relativistic and bound state corrections in the heavy quark $`b`$, $`c`$ fragmentation processes. The amplitude of heavy quark fragmentation is obtained in a new form (13) which accounts all possible relativistic factors for the calculation of relativistic corrections to the fragmentation functions. Let us summarize several peculiarities related to the calculation performed above.
1. We obtain the heavy quark fragmentation functions for both heavy quarks $`b`$ and $`c`$ which fragment to pseudoscalar and vector heavy mesons starting with the meson production amplitude (1).
2. All possible sources of relativistic corrections including the transformation factors for the two quark bound state wave function have been taken into account.
3. We investigated the role of relativistic effects in the fragmentation probabilities over two variables: the longitudinal momentum $`z`$ and transverse momentum $`p_T`$ of the heavy meson.
Analyzing the obtained analytical expressions for the fragmentation functions both in longitudinal and transverse momentum we can point out that the calculated corrections for all vector and pseudoscalar mesons are not exceeding 20 $`\%`$ the leading order contribution. The numerical value of the binding energy correction is dependent on the initial choice of the heavy quark masses because $`W=(Mm_1m_2)`$. So, for example, the binding energy corrections are extremely small for the charmonium production in our model where $`W_{J/\mathrm{\Psi }}=0.003`$ Gev. The relativistic correction is essentially more important in this case. Earlier the binding energy and relativistic corrections were studied in Ref.Bashir for the fragmentation functions of a charm quark to decay into $`\eta _c`$ and $`J/\mathrm{\Psi }`$. The comparison of our results with the calculation in Ref.Bashir shows that in the case of the reaction $`\overline{c}J/\mathrm{\Psi }`$ the binding corrections are numerically close, but relativistic corrections connected with the expression (60) are essentially different both in the sign and numerical value. Our relativistic correction coincides in the sign with the leading order contribution and is numerically three times smaller than in Ref.Bashir . Moreover, our numerical estimations of analytical relations are based on a different numerical value for the expression (60). In the paper Bashir the parameter $`𝐩^2`$ was fixed (using the mass of $`c`$-quark $`m_c=1.43GeV`$) by means of a condition analogous to our equation (60) without the addendum proportional to the parameter $`B`$ entering in the confinement part of the potential. The results obtained in the present study evidently show that the relativistic plus bound state corrections lead to the systematic increase of the fragmentation probabilities for the pseudoscalar and vector mesons. Our total fragmentation functions for the decays $`\overline{c}(c\overline{c})`$, $`\overline{b}(\overline{b}c)`$, $`\overline{b}(b\overline{b})`$ retain the initial shape of the leading order contribution. In the production of vector mesons both corrections proportional to $`W`$ and $`𝐩^2`$ have the same sign giving us a more essential modification of the leading order contribution in the comparison with the pseudoscalar meson production.
The fragmentation functions (19), (27), (42) depend not only on $`z`$ but also on the factorization scale $`\mu `$. They should be considered at a scale $`\mu `$ of the order of the heavy quark masses. The evolution of the fragmentation functions to the scale $`\mu =M_Z/2`$ is determined by the DGLAP equation DGLAP :
$$\mu ^2\frac{}{\mu ^2}D_{QH}(z,\mu ^2)=_z^1\frac{dy}{y}P_{QQ}(\frac{z}{y},\mu )D_{QH}(y,\mu ^2),$$
(63)
where $`P_{QQ}(x)`$ is the quark splitting function LP . The modification of the z-shape of the fragmentation functions is shown in Fig.6. The average values of the momentum fraction for the production of different heavy mesons at the scale $`\mu =M_Z/2`$ are the following: $`<z>=0.50`$ $`(J/\mathrm{\Psi })`$, $`<z>=0.46`$ $`(\eta _c)`$, $`<z>=0.63`$ $`(B_c^{})`$, $`<z>=0.59`$ $`(B_c)`$, $`<z>=0.56`$ $`(\mathrm{{\rm Y}})`$, $`<z>=0.52`$ $`(\eta _b)`$.
In the case of $`(\overline{c}c)`$ or $`(\overline{b}b)`$ mesons Eqs.(29) and (46) acquire a more simple form:
$$\mathrm{\Omega }_V=_0^1D_{\overline{Q}V(\overline{Q}Q)}(z)dz=\frac{32\alpha _s^2|\mathrm{\Psi }(0)|^2}{27m_2^3}[\frac{1189}{30}57\mathrm{ln}2+$$
(64)
$$+\frac{W}{m_2}(134\mathrm{ln}2\frac{78149}{840})+\frac{𝐩^2}{m_2^2}(\frac{1078}{9}\mathrm{ln}2\frac{78416}{945})].$$
$$\mathrm{\Omega }_P=_0^1D_{\overline{Q}P(\overline{Q}Q)}(z)dz=\frac{8\alpha _s^2|\mathrm{\Psi }(0)|^2}{81m_2^3}[\frac{1546}{5}444\mathrm{ln}2+$$
(65)
$$+\frac{W}{m_2}(104\mathrm{ln}2+\frac{15581}{240})+\frac{𝐩^2}{m_2^2}(\frac{16}{9}\mathrm{ln}2\frac{139}{48})].$$
Numerical values of the total fragmentation probabilities are presented in Table I. The evolution conserves the integral probabilities $`\mathrm{\Omega }_V`$ and $`\mathrm{\Omega }_P`$ of the fragmentation. Using expressions (29) and (46) we obtain the significant experimental ratio
$$\eta (r)=\frac{\mathrm{\Omega }_V}{\mathrm{\Omega }_V+\mathrm{\Omega }_P},$$
(66)
which predicts the relative number of vector and pseudoscalar mesons. For the $`(\overline{c}c)`$, $`(\overline{b}c)`$ and $`(\overline{b}b)`$ mesons the ratio (64) gives the following numbers: 0.46, 0.58, 0.52. The obtained results for the relativistic and bound state corrections to the different heavy quark fragmentation functions are shown in Figs.2-5. Relative order contributions of relativistic and binding corrections are the biggest for the $`(\overline{b}c)`$ mesons because of the little growth the parameter (60) and bound state energy $`W`$ as compared to $`(c\overline{c})`$ states. The decrease of these corrections in the bottomonium is explained by the increase of heavy quark mass ($`m_cm_b`$). All considered effects in the production of $`(\overline{b}c)`$, $`(\overline{b}b)`$ and $`(\overline{c}c)`$ mesons are computed to the leading order in $`\alpha _s`$ in the color singlet model. For S-states like $`J/\mathrm{\Psi }`$, $`\eta _c`$, $`B_c`$, $`B_c^{}`$, $`\eta _b`$ and $`\mathrm{{\rm Y}}`$ the color-octet terms are suppressed relative to the color singlet terms by a factor of the second order over the relative velocity $`v_Q`$ BFY . Our results should be useful for the comparison with more accurate $`(\overline{c}c)`$, $`(\overline{b}b)`$ and $`(\overline{b}c)`$ meson production measurements in $`Z^0`$ decays or in $`p\overline{p}`$ collisions at the Tevatron.
###### Acknowledgements.
I am grateful to D.Ebert, R.N.Faustov, V.O.Galkin for a careful reading of the manuscript, useful remarks and suggestions, and to I.B.Khriplovich, V.V.Kiselev, A.K.Likhoded, V.A.Saleev for useful discussions of different questions regarding this study. The author thanks the colleagues from the Institute of Physics of the Humboldt University in Berlin for warm hospitality. The work is performed under the financial support of the Deutsche Forschungsgemeinschaft under contract Eb 139/2-3.
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# 𝐵-meson distribution amplitudesLectures given at the Int. School on Heavy Quark Physics, Dubna, June 2005
## 1 Introduction
Many exclusive $`B`$-decay amplitudes in the framework of SCET contain $`B`$-meson distribution amplitudes . The amplitude of the decay $`B\gamma l\overline{\nu }`$ at large photon energies is given, up to power corrections, by a convolution of a hard part (perturbatively calculable) and the $`B`$-meson distribution amplitude . Amplitudes of some decays, e.g., $`B\pi l\overline{\nu }`$ at large pion energies, contain both factorizable and non-factorizable contributions . Factorizable parts of decay amplitudes contain light-cone distribution amplitudes of the initial $`B`$-meson and final hadron(s). They describe large-distance (soft) structure of these hadrons, and cannot be calculated in perturbation theory. The theory of hadronic distribution amplitudes in QCD is reviewed in .
Quark–antiquark distribution amplitudes of $`B`$-meson in HQET were introduced and investigated in . They are defined as Fourier transforms of matrix elements of some gauge-invariant bilocal operators between $`B`$-meson and vacuum. Renormalization of these operators was calculated in one-loop approximation. However, an unusual term $`1/\epsilon ^2`$ in the one-loop renormalization constants was erroneously omitted in . The correct evolution equation for the leading-twist distribution amplitude was derived in at one loop. The evolution kernel contains, in addition to terms obtained earlier , an unusual term $`\mathrm{log}(\omega /\mu )\delta (\omega \omega ^{})`$. The method of solution of the evolution equation is also discussed in in detail.
Quark–antiquark–gluon distribution amplitudes of $`B`$-meson and their relations to quark–antiquark ones (based on equations of motion) are discussed in .
Sum rules for the quark–antiquark distribution amplitudes were obtained in . A simple model of these distribution amplitudes at a low normalization scale (of order of hadronic scale) was proposed. Radiative corrections to the perturbative term and the quark-condensate term were later calculated .
In these lectures, we first briefly discuss what is HQET<sup>1</sup><sup>1</sup>1Unlike most texts on HQET, we consider not a heavy quark but a heavy antiquark. Of course, this makes no difference, but the active participant in the distribution amplitudes is the light quark (and, possibly, a gluon etc.), and this choice makes notation more natural. (Sect. 2). A much more detailed presentation can be found in the textbooks . After a short discussion of $`f_B`$ (Sect 3), quark–antiquark distribution amplitudes are introduced (Sect. 4). Quark–antiquark–gluon distribution amplitudes and their relations to two-particle ones are discussed in Sect. 5. Sect. 6 is the central (and longest) one. Here renormalization of light-cone bilocal quark operators in HQET is considered in one-loop approximation. A detailed derivation of the evolution kernel is presented, based on the methods of . Finally, sum rules for the distribution amplitudes are briefly discussed in Sect. 7.
## 2 Heavy Quark Effective Theory
Let’s consider a heavy antiquark with mass $`m`$ and momentum
$$p=mv+k$$
(2.1)
in QCD. Here $`v`$ is a fixed 4-velocity ($`v^2=1`$), and the residual momentum $`k`$, as well as momenta of all light quarks and gluons, are supposed to be small (compared to $`m`$). The propagator of this heavy antiquark can be written as
$$\frac{mm\text{/}v\text{/}k}{(mv+k)^2m^2+i0}=\frac{1\text{/}v}{2}\frac{1}{kv+i0}+𝒪\left(\frac{1}{m}\right).$$
(2.2)
The leading term here is the HQET propagator. This can be graphically presented as<sup>2</sup><sup>2</sup>2The arrow here is somewhat misleading: there is a *particle* called antiquark and propagating from left to right; it has no antiparticle.
$$\text{}=\text{}+𝒪\left(\frac{1}{m}\right).$$
(2.3)
In a vertex sandwiched between such propagators, we may substitute
$$\frac{1\text{/}v}{2}\gamma ^\mu \frac{1\text{/}v}{2}=\frac{1\text{/}v}{2}(v^\mu )\frac{1\text{/}v}{2}.$$
(2.4)
Projector can also be inserted near external legs, so that all QCD vertices can be replaced by the HQET ones,
$$\text{}=ig_0t^a(v^\mu ),$$
(2.5)
up to $`𝒪(k/m)`$ corrections.
These Feynman rules can be obtained from the Lagrangian
$$L=\overline{Q}_viv\stackrel{}{D}Q_v+(\text{light fields})$$
(2.6)
The antiquark field is our main field here, it satisfies
$$\overline{Q}_v\text{/}v=\overline{Q}_v;$$
(2.7)
$`Q_v`$ is the conjugate field for $`\overline{Q}_v`$. Here the covariant derivatives are
$$D_\mu q=\left(_\mu iA_\mu \right)q,\overline{q}\stackrel{}{D}_\mu =\overline{q}\left(\stackrel{}{}_\mu +iA_\mu \right),A_\mu =g_0A_{0\mu }^at^a.$$
(2.8)
QCD tree diagrams are reproduced by HQET up to $`𝒪(k_i/m)`$ corrections (they can also be reproduced, if we add the appropriate $`1/m`$ terms into the Lagrangian).
The heavy quark chromomagnetic moment is $`1/m`$ by dimensionality. At the leading order in $`1/m`$, the heavy-quark spin does not interact with the gluon field. Therefore, it may be rotated at will, without changing the physics (heavy-quark spin symmetry). It may even be switched off (superflavour symmetry). We shall work with the spinless heavy antiquark,
$$L=Q_v^{}iv\stackrel{}{D}Q_v+(\text{light fields}),$$
(2.9)
during most of these lectures, because this greatly simplifies reasoning and calculations. Here again $`Q_v^{}`$ is the main (scalar) field, and $`Q_v`$ is its conjugate.
So, tree QCD diagrams, expanded in $`k_i/m`$ to some order, are reproduced by the corresponding HQET diagrams. But what about loops? Here things are not so simple . Let’s consider, for example, the heavy–light two-point diagram (Fig. 1) with $`p=mv+k`$, where the residual momentum $`k`$ is small. By choosing $`v`$ along $`p`$ we can always ensure $`k=\omega v`$. Let’s consider the integral
$$I=\frac{im^2}{\pi ^{d/2}}\frac{d^dl}{\left[m^2(mv+k+l)^2\right]^2(l^2)},$$
(2.10)
which has neither ultraviolet (UV) nor infrared (IR) divergences. There are two regions of the loop momentum $`l`$ in this integral:
* Hard region $`lm`$. Expanding the integrand in $`km`$, $`l`$, we have
$$\frac{1}{\left[m^2(mv+l)^2\right]^2(l^2)}+4\frac{(m+lv)\omega }{\left[m^2(mv+l)^2\right]^3(l^2)}+\mathrm{}$$
(2.11)
* Soft region $`l\omega `$. Expanding the integrand in $`k`$, $`lm`$, we have
$$\frac{1}{\left[2m(k+l)v\right]^2(l^2)}+2\frac{(k+l)^2}{\left[2m(k+l)v\right]^3(l^2)}+\mathrm{}$$
(2.12)
The contribution of the hard region is
$$\begin{array}{cc}& I_h=m^{2\epsilon }\left[M(2,1)+2\frac{\omega }{m}\left[M(3,0)M(2,1)+2M(3,1)\right]+\mathrm{}\right],\hfill \\ & \frac{\mu ^{2\epsilon }I_h}{\mathrm{\Gamma }(1+\epsilon )}=\frac{1}{2\epsilon }+\mathrm{log}\frac{m}{\mu }+\left(\frac{1}{\epsilon }2\mathrm{log}\frac{m}{\mu }1\right)\frac{\omega }{m}+\mathrm{}\hfill \end{array}$$
(2.13)
where the on-shell massive two-point integrals (Fig. 2)
$$\begin{array}{cc}& \frac{d^dl}{D_1^{n_1}D_2^{n_2}}=i\pi ^{d/2}m^{d2(n_1+n_2)}M(n_1,n_2),\hfill \\ & D_1=m^2(l+mv)^2i0,D_2=l^2i0\hfill \end{array}$$
(2.14)
are
$$M(n_1,n_2)=\frac{\mathrm{\Gamma }(dn_12n2)\mathrm{\Gamma }(d/2+n_1+n_2)}{\mathrm{\Gamma }(n_1)\mathrm{\Gamma }(dn_1n_2)}.$$
(2.15)
This contribution is IR divergent.
The contribution of the soft region is
$$\begin{array}{cc}& I_s=(2\omega )^{2\epsilon }\left[I(2,1)+\frac{\omega }{m}\left[I(3,1)2I(2,1)\right]+\mathrm{}\right],\hfill \\ & \frac{\mu ^{2\epsilon }I_s}{\mathrm{\Gamma }(1+\epsilon )}=\frac{1}{2\epsilon }\mathrm{log}\frac{2\omega }{\mu }\left(\frac{1}{\epsilon }2\mathrm{log}\frac{2\omega }{\mu }\frac{1}{2}\right)\frac{\omega }{m}+\mathrm{}\hfill \end{array}$$
(2.16)
where the HQET two-point integrals (Fig. 3, $`\omega =kv`$)
$$\begin{array}{cc}& \frac{d^dl}{D_1^{n_1}D_2^{n_2}}=i\pi ^{d/2}(2\omega )^{dn_12n_2}I(n_1,n_2),\hfill \\ & D_1=2(l+p)vi0,D_2=l^2i0\hfill \end{array}$$
(2.17)
are
$$I(n_1,n_2)=\frac{\mathrm{\Gamma }(d+n_1+2n_2)\mathrm{\Gamma }(d/2n_2)}{\mathrm{\Gamma }(n_1)\mathrm{\Gamma }(n_2)}.$$
(2.18)
This contribution is UV divergent.
The complete result is finite:
$$I=\mathrm{log}\frac{2\omega }{m}+\left(2\mathrm{log}\frac{2\omega }{m}\frac{1}{2}\right)\frac{\omega }{m}+\mathrm{}$$
(2.19)
This expansion can be easily continued if desired.
What about higher loops? Let’s consider, for example, the two-loop two-point diagram with a small external residual momentum. There are several regions of the loop momenta in this diagram (Fig. 4). In each of them, some momenta are hard ($`m`$), some are soft ($`\omega `$). As we have already seen, a heavy-quark line with a soft residual momentum becomes an HQET line. Soft massless lines are shown by dashed lines in the figure.
For example, in the second diagram in Fig. 4, the left loop is hard. It contains a single scale $`m`$. All external momenta of this loop are soft, including those of the lines belonging to the other loop. We can expand the integrand in these small momenta; after taking the loop integral, we obtain a polynomial in these momenta. Now from the point of view of the soft loop (large distances) this hard loop is just a local vertex. This loop contains a single scale $`\omega `$, and we can calculate it, obtaining a non-analytical function of the external residual energy $`\omega `$.
In a general multiloop diagram, hard lines must always form loops, so that momentum conservation can hold after neglecting all soft momenta. There can be several disconnected hard parts; each of them must contain at least one heavy line (a subdiagram consisting of only light lines and having soft external momenta has no reason to be hard). From the point of view of the soft part, hard parts are just local vertices. There may appear a soft subdiagram connected to the rest of the diagram only at such a vertex. Such a subdiagram is scaleless and hence vanish. For example, we could consider one more region in Fig. 4: when the only soft line is the middle light line, all the rest are hard. But this soft line forms a loop containing one local vertex (the integrated hard loop), and hence it vanishes.
In the usual HQET formalism, Lagrangian contains local operators multiplied by matching coefficients. QCD operators are also expanded in HQET operators with matching coefficients. These matching coefficients are the only quantities in the theory which depend on the hard scale $`m`$. Diagrammatically, they come from hard loops in QCD diagrams. Local operators produce vertices polynomial in their external momenta. They appear in HQET diagrams, which contain only the soft scale $`\omega `$. These HQET diagrams are soft parts of QCD diagrams.
## 3 $`B`$-meson decay constant
During these lectures, we shall mostly live in a world with a heavy antiquark having $`j^P=0^+`$. Physics in the real world is the same, up to $`1/m`$ corrections. We shall work in the $`v`$ rest frame.
The ground-state $`S`$-wave $`\overline{Q}q`$ meson has the quantum numbers $`j^P=\frac{1}{2}^+`$. There are 2 $`P`$-wave excited mesons with $`j^P=\frac{1}{2}^{}`$ and $`\frac{3}{2}^{}`$. The heavy–light quark current
$$j=Q_v^{}q$$
(3.1)
has no definite parity; the currents with parity $`P=\pm 1`$ are
$$j_P=\frac{1+P\gamma ^0}{2}j.$$
(3.2)
They have the quantum numbers of $`S`$-wave $`\frac{1}{2}^+`$ mesons and $`P`$-wave $`\frac{1}{2}^{}`$ mesons. The ground-state meson $`M`$ has a Dirac wave function $`u`$ which satisfies $`\gamma ^0u=u`$ and is normalized by $`\overline{u}u=1`$. The matrix element of $`j`$ from $`M`$ to vacuum is
$$<0|j|M>=Fu,$$
(3.3)
where the one-meson states are normalized by the non-relativistic condition
$$<M,\stackrel{}{p}^{}|M,\stackrel{}{p}>=(2\pi )^3\delta (\stackrel{}{p}^{}\stackrel{}{p}).$$
(3.4)
The correlator of the heavy–light currents (Fig. 5),
$$i<Tj(x)\overline{ȷ}(0)>=\delta (\stackrel{}{x})\mathrm{\Pi }(x^0),$$
(3.5)
has the structure
$$\mathrm{\Pi }(x^0)=A+B\text{/}v.$$
(3.6)
Therefore, the correlators of the currents with definite parity are
$$i<Tj_P(x)\overline{ȷ}_P(0)>=\delta (\stackrel{}{x})\mathrm{\Pi }_P(x^0),\mathrm{\Pi }_P=A+PB=\frac{1}{4}\mathrm{Tr}(1+P\gamma ^0)\mathrm{\Pi }.$$
(3.7)
The spectral density of, say, the correlator with $`P=+1`$ is
$$\rho _+(\epsilon )=F^2\delta (\epsilon \overline{\mathrm{\Lambda }})+\mathrm{}$$
(3.8)
where $`\overline{\mathrm{\Lambda }}`$ is the residual energy of the ground-state meson, and the dots mean contribution of excited states.
Now we shall return to the real world with $`j^P=\frac{1}{2}^{}`$ heavy antiquark for a while, but still with $`m=\mathrm{}`$. The $`S`$-wave ground-state meson turns into a degenerate doublet with $`j^P=0^{}`$, $`1^{}`$. The two $`P`$-wave mesons turn into two degenerate doublets: with $`j^P=0^+`$, $`1^+`$ and with $`j^P=1^+`$, $`2^+`$. All heavy–light currents $`\overline{Q}_v\mathrm{\Gamma }q`$ reduce, due to (2.7), to 4 ones with
$$\mathrm{\Gamma }=\gamma _5,\stackrel{}{\gamma }\text{and}\mathrm{\Gamma }=1,\stackrel{}{\gamma }\gamma _5.$$
(3.9)
The first 2 currents, with $`\mathrm{\Gamma }`$ anticommuting with $`\gamma ^0`$, have the quantum numbers of the ground-state $`0^{}`$, $`1^{}`$ mesons; the second 2 currents, with $`\mathrm{\Gamma }`$ commuting with $`\gamma ^0`$, have the quantum numbers of the $`P`$-wave $`0^+`$, $`1^+`$ mesons.
The correlators (Fig. 5) are
$$i<Tj_2(x)j_1^+(0)>=\delta (\stackrel{}{x})\mathrm{\Pi }_{12}(x^0),\mathrm{\Pi }_{12}=\mathrm{Tr}\overline{\mathrm{\Gamma }}_1\frac{1\gamma ^0}{2}\mathrm{\Gamma }_2\mathrm{\Pi },$$
(3.10)
or
$$\mathrm{\Pi }_{12}=\mathrm{\Pi }_P\mathrm{Tr}\overline{\mathrm{\Gamma }}_1\frac{1\gamma ^0}{2}\mathrm{\Gamma }_2,$$
(3.11)
where $`P=+1`$ for $`\mathrm{\Gamma }`$ anticommuting with $`\gamma ^0`$ and $`1`$ for commuting. For $`\mathrm{\Gamma }=\gamma _5`$ and $`\gamma ^i`$, the correlators are $`2\mathrm{\Pi }_+`$ and $`2\mathrm{\Pi }_+\delta ^{ij}`$. Their spectral densities are $`F_B^2\delta (\epsilon \overline{\mathrm{\Lambda }})+\mathrm{}`$ and $`F_B^{}^2\delta (\epsilon \overline{\mathrm{\Lambda }})\delta ^{ij}+\mathrm{}`$, where
$$<0|\overline{Q}_v\gamma _5q|B>=F_B,<0|\overline{Q}_v\stackrel{}{\gamma }q|B^{}>=F_B^{}\stackrel{}{e},$$
(3.12)
and dots mean contribution of higher states. Therefore,
$$F_B=F_B^{}=\sqrt{2}F.$$
(3.13)
The usual definition of the decay constants is
$$<0|\overline{Q}_v\gamma ^\mu \gamma _5q|B>_r=if_Bp^\mu ,<0|\overline{Q}_v\gamma ^\mu q|B^{}>_r=imf_B^{}e^\mu ,$$
(3.14)
where the relativistic normalization
$${}_{r}{}^{}<B,p^{}|B,p>_r=(2\pi )^32p^0\delta (\stackrel{}{p}^{}\stackrel{}{p})$$
(3.15)
of single-meson states is used (this normalization becomes meaningless in the limit $`m\mathrm{}`$, and therefore cannot be used in HQET). Comparing (3.14) for $`B`$ or $`B^{}`$ at rest with (3.12), we arrive at
$$f_B=f_B^{}=\frac{2F}{\sqrt{m}},$$
(3.16)
up to $`1/m`$ corrections.
Let’s consider the correlator (3.11) with $`P=+1`$. Equating the ground-state contributions to the spectral densities of the left-hand side and the right-hand one, we obtain
$$<0|j_2|M><M|j_1^+|0>=F^2\mathrm{Tr}\mathrm{\Gamma }_2\frac{1+\gamma ^0}{2}\overline{\mathrm{\Gamma }}_1.$$
(3.17)
Therefore,
$$<0|\overline{Q}_v\mathrm{\Gamma }q|M>=\frac{F}{\sqrt{2}}\mathrm{Tr}\mathrm{\Gamma }\frac{1+\gamma ^0}{2}\mathrm{\Gamma }_M,$$
(3.18)
where the matrix
$$\mathrm{\Gamma }_M=\{\begin{array}{cc}i\gamma _5\hfill & \text{for }B\hfill \\ i\text{/}e\hfill & \text{for }B^{}\hfill \end{array}$$
(3.19)
is defined up to a phase factor. We can re-write this result for the relativistic normalization of the meson state, in covariant notation:
$$<0|\overline{Q}_v\mathrm{\Gamma }q|M>_r=\sqrt{m}F\mathrm{Tr}\mathrm{\Gamma },$$
(3.20)
where
$$=\frac{1+\text{/}v}{2}\{\begin{array}{cc}i\gamma _5\hfill & \text{for }B\hfill \\ i\text{/}e\hfill & \text{for }B^{}\hfill \end{array}$$
(3.21)
(of course, one can re-define the phases of $`|M>`$ and hence of $``$).
## 4 Quark–antiquark distribution amplitudes
After safely returning to the ideal world with a $`0^+`$ heavy antiquark, we want to invent an operator which probes more details of the structure of $`B`$-meson than the local current (3.1). To this end, we consider a bilocal gauge-invariant operator
$$\stackrel{~}{O}(t)=Q_v^{}(0)[0,z]q(z),$$
(4.1)
where
$$[x,y]=P\mathrm{exp}\left[i_x^yA_\mu (z)𝑑z^\mu \right],$$
(4.2)
and
$$z^2=0,t=vz.$$
(4.3)
Its matrix element from the ground-state meson to vacuum has 2 Dirac structures:
$$<0|\stackrel{~}{O}(t)|M>=F\left[\stackrel{~}{\phi }_+(t)+\frac{\stackrel{~}{\phi }_{}(t)\stackrel{~}{\phi }_+(t)}{2t}\text{/}z\right]u.$$
(4.4)
because $`\text{/}vu=u`$.
In what follows, we shall often use light-front components of vectors. Let’s introduce (in the $`v`$ rest frame) two light-like vectors
$$\begin{array}{cc}& n_\pm ^\mu =(1,1,\stackrel{}{0}),\hfill \\ & n_+^2=n_{}^2=0,n_+n_{}=2.\hfill \end{array}$$
(4.5)
Light-front components of any vector $`a`$ are defined as
$$a_\pm =an_\pm =a^0\pm a^1.$$
(4.6)
We have
$$\begin{array}{cc}& a^\mu =\frac{1}{2}\left(a_+n_{}^\mu +a_{}n_+^\mu \right)+a_{}^\mu ,\hfill \\ & ab=\frac{1}{2}\left(a_+b_{}+a_{}b_+\right)\stackrel{}{a}_{}\stackrel{}{b}_{}.\hfill \end{array}$$
(4.7)
In particular,
$$v^\mu =\frac{1}{2}\left(n_+^\mu +n_{}^\mu \right),v_+=v_{}=1,\stackrel{}{v}_{}=\stackrel{}{0}.$$
(4.8)
We shall also use light-front components of $`\gamma ^\mu `$:
$$\gamma _\pm =\gamma n_\pm =\text{/}n_\pm .$$
(4.9)
The definition (4.4) can be re-written as
$$<0|\stackrel{~}{O}(t)|M>=\frac{1}{2}F\left[\stackrel{~}{\phi }_+(t)\gamma _{}+\stackrel{~}{\phi }_{}(t)\gamma _+\right]u.$$
(4.10)
If we introduce the operators
$$\stackrel{~}{O}_\pm (t)=\gamma _\pm \stackrel{~}{O}(t),$$
(4.11)
then
$$<0|\stackrel{~}{O}_\pm (t)|M>=F\stackrel{~}{\phi }_\pm (t)\gamma _\pm u.$$
(4.12)
The $`B`$-meson distribution amplitudes are the Fourier transforms of these functions:
$$\phi _\pm (\omega )=\frac{1}{2\pi }\stackrel{~}{\phi }_\pm (t)e^{i\omega t}𝑑t,\stackrel{~}{\phi }_\pm (t)=\phi _\pm (\omega )e^{i\omega t}𝑑\omega .$$
(4.13)
They are normalized by
$$\stackrel{~}{\phi }_\pm (0)=_0^{\mathrm{}}\phi _\pm (\omega )𝑑\omega =1.$$
(4.14)
The function $`\phi _+(\omega )`$ is the leading-twist distribution amplitude, and $`\phi _{}(\omega )`$ – the subleading-twist one (though there is no good definition of twist in HQET). We can formally introduce the operators
$$\begin{array}{cc}& O_\pm (\omega )=\frac{1}{2\pi }\stackrel{~}{O}_\pm (t)e^{i\omega t}𝑑t=Q_v^{}(0)\gamma _\pm \delta (iD_+\omega )q(0),\hfill \\ & \stackrel{~}{O}_\pm (t)=O_\pm (\omega )e^{i\omega t}𝑑\omega ,\hfill \end{array}$$
(4.15)
then
$$<0|O_\pm (\omega )|M>=F\phi _\pm (\omega )\gamma _\pm u.$$
(4.16)
The distribution amplitudes describe the distribution in the light-front component $`p_+`$ of the light-quark momentum in $`B`$-meson.
The expansion of the operators (4.11) in $`t`$ reads
$$\begin{array}{cc}& \stackrel{~}{O}_\pm (t)=\underset{n=0}{\overset{\mathrm{}}{}}O_\pm ^{(n)}\frac{(it)^n}{n!},\hfill \\ & O_\pm ^{(n)}=O_\pm (\omega )\omega ^n𝑑\omega =Q_v^{}\gamma _\pm (iD_+)^nq,\hfill \end{array}$$
(4.17)
or for matrix elements
$$\begin{array}{cc}& \stackrel{~}{\phi }_\pm (t)=\underset{n=0}{\overset{\mathrm{}}{}}<\omega ^n>_\pm \frac{(it)^n}{n!},\hfill \\ & <\omega ^n>_\pm =_0^{\mathrm{}}\phi _\pm (\omega )\omega ^n𝑑\omega ,\hfill \\ & <0|O_\pm ^{(n)}|M>=F<\omega ^n>_\pm \gamma _\pm u.\hfill \end{array}$$
(4.18)
We can also reconstruct $`O_\pm (\omega )`$ from $`O_\pm ^{(n)}`$:
$$O_\pm (\omega )=_i\mathrm{}^{+i\mathrm{}}O_\pm ^{(n)}\omega ^{n1}\frac{dn}{2\pi i},$$
(4.19)
or for matrix elements
$$\phi _\pm (\omega )=_i\mathrm{}^{+i\mathrm{}}<\omega ^n>_\pm \omega ^{n1}\frac{dn}{2\pi i}$$
(4.20)
(you can easily check this by substituting (4.17) into (4.19); integration in $`dn`$ yields $`\delta (\mathrm{log}(\omega ^{}/\omega ))`$).
The first moments $`<\omega >_\pm `$ can be found from the equations of motion. The equation of motion for the heavy antiquark is
$$Q_v^{}\stackrel{}{D}_0=0,$$
(4.21)
and therefore we obtain
$$<0|Q_v^{}D_0q|M>=<0|_0(Q_v^{}q)|M>=iF\overline{\mathrm{\Lambda }}u.$$
(4.22)
The vector part has the structure
$$<0|Q_v^{}\stackrel{}{D}q|M>=aF\stackrel{}{\gamma }u.$$
(4.23)
The equation of motion of the light quark is
$$<0|Q_v^{}\text{ /}Dq|M>=0,$$
(4.24)
and using $`\text{ /}D=D_0\gamma ^0\stackrel{}{D}\stackrel{}{\gamma }`$ we obtain
$$<0|Q_v^{}\stackrel{}{D}q|M>=\frac{i}{3}F\overline{\mathrm{\Lambda }}\stackrel{}{\gamma }u.$$
(4.25)
Finally, we arrive at
$$<\omega >_+=\frac{4}{3}\overline{\mathrm{\Lambda }},<\omega >_{}=\frac{2}{3}\overline{\mathrm{\Lambda }}.$$
(4.26)
Let’s now consider the second moments. They involve two new (nonperturbative) hadronic parameters:
$$<0|Q_v^{}\stackrel{}{E}\stackrel{}{\alpha }q|M>=iF\lambda _E^2u,<0|Q_v^{}\stackrel{}{H}\stackrel{}{\sigma }q|M>=F\lambda _H^2u,$$
(4.27)
where
$$\stackrel{}{E}=i[D_0,\stackrel{}{D}],\stackrel{}{H}=i\stackrel{}{D}\times \stackrel{}{D},\stackrel{}{\alpha }=\gamma ^0\stackrel{}{\gamma },\stackrel{}{\sigma }=\stackrel{}{\gamma }\gamma _5\gamma ^0.$$
(4.28)
Now we can calculate all matrix elements with 2 derivatives. From (4.22) and (4.25), using the heavy-antiquark equation of motion (4.21), we immediately find
$$<0|Q_v^{}D_0^2q|M>=F\overline{\mathrm{\Lambda }}^2u,<0|Q_v^{}D_0\stackrel{}{D}q|M>=\frac{1}{3}F\overline{\mathrm{\Lambda }}^2\stackrel{}{\gamma }u.$$
(4.29)
Using the definition (4.27) of $`\lambda _E^2`$, we immediately find
$$<0|Q_v^{}\stackrel{}{D}D_0q|M>=\frac{1}{3}F\left(\overline{\mathrm{\Lambda }}^2+\lambda _E^2\right)\stackrel{}{\gamma }u.$$
(4.30)
The second spatial derivatives have the structure
$$<0|Q_v^{}D^iD^jq|M>=F\left(b\delta ^{ij}\frac{i}{6}\lambda _H^2\epsilon ^{ijk}\sigma ^k\right)u,$$
(4.31)
where the second coefficient follows from the definition (4.27) of $`\lambda _H^2`$. We find $`b`$ using the light-quark equation of motion (4.24):
$$<0|Q_v^{}D^iD^jq|M>=\frac{1}{3}F\left[\left(\overline{\mathrm{\Lambda }}^2+\lambda _E^2+\lambda _H^2\right)\delta ^{ij}+\frac{i}{2}\lambda _H^2\epsilon ^{ijk}\sigma ^k\right]u.$$
(4.32)
Finally, we arrive at
$$<\omega ^2>_+=2\overline{\mathrm{\Lambda }}^2+\frac{2}{3}\lambda _E^2+\frac{1}{3}\lambda _H^2,<\omega ^2>_{}=\frac{2}{3}\overline{\mathrm{\Lambda }}^2+\frac{1}{3}\lambda _H^2.$$
(4.33)
What about $`B`$-meson distribution amplitudes in the real world with a $`\frac{1}{2}^{}`$ heavy antiquark? It has 4 distribution amplitudes, as any pseudoscalar meson:
$$\begin{array}{cc}& <0|Q_v^{}(0)[0,z]\gamma _5q(z)|B>_r=if_Bm\stackrel{~}{\phi }_P,\hfill \\ & <0|Q_v^{}(0)[0,z]\gamma ^\mu \gamma _5q(z)|B>_r=f_B\left[i\stackrel{~}{\phi }_{A1}p^\mu m\stackrel{~}{\phi }_{A2}z^\mu \right],\hfill \\ & <0|Q_v^{}(0)[0,z]\sigma ^{\mu \nu }\gamma _5q(z)|B>_r=if_B\stackrel{~}{\phi }_T(p^\mu z^\nu p^\nu z^\mu ).\hfill \end{array}$$
(4.34)
Similarly to (3.20), we have
$$<0|Q_v^{}(0)[0,z]\mathrm{\Gamma }q(z)|M>_r=F\mathrm{Tr}\mathrm{\Gamma }\left[\stackrel{~}{\phi }_++\frac{\stackrel{~}{\phi }_{}\stackrel{~}{\phi }_+}{2t}\text{/}z\right].$$
(4.35)
Therefore these 4 QCD distribution amplitudes can be expressed via 2 HQET ones:
$$\begin{array}{cc}& \stackrel{~}{\phi }_P=\frac{\stackrel{~}{\phi }_+(t)+\stackrel{~}{\phi }_{}(t)}{2},\stackrel{~}{\phi }_{A1}=\stackrel{~}{\phi }_+(t),\hfill \\ & \stackrel{~}{\phi }_{A2}=\stackrel{~}{\phi }_T=\frac{i}{2}\frac{\stackrel{~}{\phi }_+(t)\stackrel{~}{\phi }_{}(t)}{t}.\hfill \end{array}$$
(4.36)
The QCD distribution amplitudes are usually considered as functions of the longitudinal momentum fraction $`x=\omega /m`$. It is usually assumed that $`\phi _{A1}(x)x`$ at $`x0`$, and $`\phi _P(x)\text{const}`$. Therefore, we shall assume that $`\phi _+(\omega )\omega `$ at $`\omega 0`$, and $`\phi _{}(\omega )\text{const}`$. The QCD distribution amplitudes are normalized as
$$\stackrel{~}{\phi }_P(0)=\stackrel{~}{\phi }_{A1}(0)=1,\stackrel{~}{\phi }_{A2}(0)=\stackrel{~}{\phi }_T(0)=\frac{\overline{\mathrm{\Lambda }}}{3}$$
(4.37)
(to derive the last formula, we used the $`t`$ expansion (4.18) and the first moments (4.26)).
The vector meson $`B^{}`$ is described by 6 distribution amplitudes in QCD. All 10 distribution amplitudes (4 for $`B`$ plus 6 for $`B^{}`$) are expressed via 2 HQET distribution amplitudes $`\phi _\pm (\omega )`$. This is a consequence of the heavy-quark spin symmetry. We don’t present formulae for $`B^{}`$ here; they can be found in .
## 5 Quark–antiquark–gluon distribution amplitudes
Now we shall discuss relations of quark–antiquark distribution amplitudes and quark–antiquark–gluon ones following from the equations of motion. In order to apply them, we need to differentiate with respect to the coordinates of the light quark and the heavy antiquark separately. Therefore, we go slightly off the light cone. The generalization of (4.4) to an arbitrary $`z^2`$ is
$$<0|Q_v^{}(0)q(z)|M>=F\left[\stackrel{~}{\phi }_+(t,z^2)+\frac{\stackrel{~}{\phi }_{}(t,z^2)\stackrel{~}{\phi }_+(t,z^2)}{2t}\text{/}z\right]u,$$
(5.1)
where $`t=vz`$, and the fixed-point gauge $`x_\mu A^\mu (x)=0`$ is used to simplify notation.
In order to apply the light-quark equation of motion (4.24), we apply
$$\gamma ^\mu \frac{}{z^\mu }$$
to this definition. The differentiation of the right-hand side is straightforward. In the left-hand side we write
$$<0|Q_v^{}(0)\gamma ^\mu (_\mu iA_\mu (z)+iA_\mu (z))q(z)|M>.$$
The first two terms yield zero. In the fixed-point gauge
$$A_\mu (z)=_0^1G_{\nu \mu }(uz)uz^\nu 𝑑u,$$
(5.2)
where $`G_{\mu \nu }=gG_{\mu \nu }^at^a`$.
The matrix element of the operator containing $`G_{\mu \nu }`$ contains 4 structures:
$$\begin{array}{cc}& <0|Q_v^{}(0)[0,uz]iG_{\nu \mu }(uz)z^\nu [uz,z]q(z)|M>=\hfill \\ & F\left[(v_\mu \text{/}zt\gamma _\mu )(\stackrel{~}{\psi }_A\stackrel{~}{\psi }_V)+i\sigma _{\mu \nu }z^\nu \stackrel{~}{\psi }_Vz_\mu \stackrel{~}{\psi }_X+\frac{z_\mu }{t}\text{/}z\stackrel{~}{\psi }_Y\right]u,\hfill \end{array}$$
(5.3)
and is parametrized by 4 quark–antiquark–gluon distribution amplitudes (note that this matrix elements vanishes when multiplied by $`z^\mu `$).
Equating the coefficients of $`u`$ and $`\text{/}zu`$, we obtain two consequences of the light-quark equation of motion:
$`\stackrel{~}{\phi }_{}^{}+{\displaystyle \frac{\stackrel{~}{\phi }_{}\stackrel{~}{\phi }_+}{t}}=2t{\displaystyle _0^1}(\stackrel{~}{\psi }_A\stackrel{~}{\psi }_V)u𝑑u,`$ (5.4)
$`\stackrel{~}{\phi }_+^{}\stackrel{~}{\phi }_{}^{}+{\displaystyle \frac{\stackrel{~}{\phi }_{}\stackrel{~}{\phi }_+}{t}}+4t{\displaystyle \frac{\stackrel{~}{\phi }_+}{z^2}}=2t{\displaystyle _0^1}(2\stackrel{~}{\psi }_V+\stackrel{~}{\psi }_A+\stackrel{~}{\psi }_X)u𝑑u.`$ (5.5)
Similarly, to use the heavy-antiquark equation of motion we apply
$$Q_v^{}(0)v^\mu \frac{}{z^\mu }q(z)=v^\mu _\mu (Q_v^{}(0)q(z))v^\mu Q_v^{}(0)(\stackrel{}{}_\mu +iA_\mu (0)iA_\mu (0))q(z).$$
The first two terms in the last operator yield zero. Now we move the heavy antiquark; the center of the fixed-point gauge must be constant, and we have to use the gauge $`(xz)^\mu A_\mu (x)=0`$:
$$A_\mu (0)=_0^1G_{\nu \mu }(uz)(1u)z^\nu 𝑑u.$$
(5.6)
We arrive at two consequences of the heavy-antiquark equation of motion:
$`\stackrel{~}{\phi }_+^{}+{\displaystyle \frac{\stackrel{~}{\phi }_{}\stackrel{~}{\phi }_+}{2t}}+i\overline{\mathrm{\Lambda }}\stackrel{~}{\phi }_++2t{\displaystyle \frac{\stackrel{~}{\phi }_+}{z^2}}=t{\displaystyle _0^1}(\stackrel{~}{\psi }_A+\stackrel{~}{\psi }_X)(1u)𝑑u,`$ (5.7)
$`\stackrel{~}{\phi }_{}^{}\stackrel{~}{\phi }_+^{}+{\displaystyle \frac{\stackrel{~}{\phi }_+\stackrel{~}{\phi }_{}}{t}}+i\overline{\mathrm{\Lambda }}(\stackrel{~}{\phi }_{}\stackrel{~}{\phi }_+)+2t\left({\displaystyle \frac{\stackrel{~}{\phi }_{}}{z^2}}{\displaystyle \frac{\stackrel{~}{\phi }_+}{z^2}}\right)=2t{\displaystyle _0^1}(\stackrel{~}{\psi }_A+\stackrel{~}{\psi }_Y)(1u)𝑑u.`$ (5.8)
The functions
$$\frac{\stackrel{~}{\phi }_\pm (t,z^2)}{z^2}|_{z^2=0}$$
are some new (non-leading) quark–antiquark distribution amplitudes. We are not interested in them here. There are 2 equations involving only our familiar distribution amplitudes $`\stackrel{~}{\phi }_\pm (t,0)`$, namely, (5.4) and a combination of (5.5) and (5.7):
$$\begin{array}{cc}& \stackrel{~}{\phi }_{}^{}+\frac{\stackrel{~}{\phi }_{}\stackrel{~}{\phi }_+}{t}=2t_0^1(\stackrel{~}{\psi }_A\stackrel{~}{\psi }_V)u𝑑u,\hfill \\ & \stackrel{~}{\phi }_+^{}+\stackrel{~}{\phi }_{}^{}+2i\overline{\mathrm{\Lambda }}\stackrel{~}{\phi }_+=2t_0^1(\stackrel{~}{\psi }_A+\stackrel{~}{\psi }_X+2\stackrel{~}{\psi }_Vu)𝑑u.\hfill \end{array}$$
(5.9)
The quark–antiquark–gluon distribution amplitudes in the momentum space are defined as
$$\stackrel{~}{\psi }_i(t,u)=\psi _i(\omega ,\xi )e^{i(\omega +\xi u)t}𝑑\omega 𝑑\xi .$$
(5.10)
Performing Fourier transform of the equations (5.9), we obtain
$$\begin{array}{cc}& \omega \frac{d\phi _{}(\omega )}{d\omega }+\phi _+(\omega )=I(\omega ),\hfill \\ & (\omega 2\overline{\mathrm{\Lambda }})\phi _+(\omega )+\omega \phi _{}(\omega )=J(\omega ),\hfill \end{array}$$
(5.11)
where
$`I(\omega )=2{\displaystyle \frac{d}{d\omega }}{\displaystyle _0^\omega }𝑑\rho {\displaystyle _{\omega \rho }^{\mathrm{}}}{\displaystyle \frac{d\xi }{\xi }}{\displaystyle \frac{}{\xi }}\left[\psi _A(\rho ,\xi )\psi _V(\rho ,\xi )\right],`$
$`J(\omega )=2{\displaystyle \frac{d}{d\omega }}{\displaystyle _0^\omega }𝑑\rho {\displaystyle _{\omega \rho }^{\mathrm{}}}{\displaystyle \frac{d\xi }{\xi }}\left[\psi _A(\rho ,\xi )+\psi _X(\rho ,\xi )\right]4{\displaystyle _0^\omega }𝑑\rho {\displaystyle _{\omega \rho }^{\mathrm{}}}{\displaystyle \frac{d\xi }{\xi }}{\displaystyle \frac{}{\xi }}\psi _V(\rho ,\xi ).`$
If we insist on having $`\phi _+(0)=0`$, then the condition
$$J(0)=2_0^{\mathrm{}}\frac{d\xi }{\xi }\left[\psi _A(0,\xi )+\psi _X(0,\xi )\right]=0$$
(5.12)
must be satisfied. It should follow from vanishing of $`\psi _{A,X}(\omega ,\xi )`$ at $`\omega 0`$ (behaviour of the three-particle distribution amplitudes at the boundaries is discussed in ).
If we new $`\psi _{A,V,X}(\omega ,\xi )`$, we could solve the equations (5.11) for $`\phi _\pm (\omega )`$. This is not very realistic, because we don’t know them. The solution is a sum of two terms:
$$\phi _\pm (\omega )=\phi _\pm ^{(WW)}(\omega )+\phi _\pm ^{(g)}(\omega ),$$
(5.13)
where $`\phi _\pm ^{(WW)}(\omega )`$ is the solution of (5.11) in the case if all quark–antiquark–gluon distribution amplitudes vanish (it is called the Wandzura–Wilczek part of the solution), and $`\phi _\pm ^{(g)}(\omega )`$ is induced by the gluonic terms. The Wandzura–Wilczek part is
$$\phi _+^{(WW)}(\omega )=\frac{\omega }{2\overline{\mathrm{\Lambda }}^2}\theta (2\overline{\mathrm{\Lambda }}\omega ),\phi _{}^{(WW)}(\omega )=\frac{2\overline{\mathrm{\Lambda }}\omega }{2\overline{\mathrm{\Lambda }}^2}\theta (2\overline{\mathrm{\Lambda }}\omega )$$
(5.14)
(Fig. 6); it satisfies the normalization conditions (4.14). Note that $`2\overline{\mathrm{\Lambda }}`$ is the maximum value of $`p_+`$ of the light quark if we assume that $`B`$-meson (having the residual energy $`\overline{\mathrm{\Lambda }}`$) consists of the on-shell heavy antiquark (always having zero residual energy) and the on-shell light quark. The gluon-induced part is given by some explicit integrals of $`\psi _{A,V,X}(\omega ,\xi )`$; we don’t present these long expressions here, they can be found in .
The moments (4.18) also consist of two contributions. The Wandzura–Wilczek parts are
$$<\omega ^n>_+^{(WW)}=\frac{2(2\overline{\mathrm{\Lambda }})^n}{n+2},<\omega ^n>_{}^{(WW)}=\frac{2(2\overline{\mathrm{\Lambda }})^n}{(n+1)(n+2)}.$$
(5.15)
The gluon-induced part does not contribute to the zeroth moments (normalization) and the first ones; its contribution to the second moments (4.33) is expressed via the normalizations of the quark–antiquark–gluon distribution amplitudes:
$$\begin{array}{cc}& \psi _A(\omega ,\xi )𝑑\omega 𝑑\xi =\frac{1}{3}\lambda _E^2,\hfill \\ & \psi _V(\omega ,\xi )𝑑\omega 𝑑\xi =\frac{1}{3}\lambda _H^2,\hfill \\ & \psi _X(\omega ,\xi )𝑑\omega 𝑑\xi =0.\hfill \end{array}$$
(5.16)
The contributions to the higher moments are expressed via the moments of $`\psi _{A,V,X}(\omega ,\xi )`$; the explicit formulae can be found in .
## 6 Evolution
Until now, we neglected renormalization of the considered operators. In this section we shall discuss renormalization of the leading-twist $`B`$-meson distribution amplitude.
Let’s summarize what we know about the bare operators of interest. There are 3 families of such operators, or 3 “representations” ($`t`$, $`\omega `$, $`n`$):
$$\begin{array}{cc}& \stackrel{~}{O}_+(t)=Q_v^{}(0)\gamma _+[0,z]q(z),\hfill \\ & O_+(\omega )=Q_v^{}(0)\gamma _+\delta (iD_+\omega )q(0),\hfill \\ & O_+^{(n)}=Q_v^{}(0)\gamma _+(iD_+)^nq(0).\hfill \end{array}$$
(6.1)
They can be converted into each other:
$$\begin{array}{cc}& \stackrel{~}{O}_+(t)=O_+(\omega )e^{i\omega t}𝑑\omega =\underset{n=0}{\overset{\mathrm{}}{}}O_+^{(n)}\frac{(it)^n}{n!},\hfill \\ & O_+(\omega )=\stackrel{~}{O}_+(t)e^{i\omega t}\frac{dt}{2\pi }=_i\mathrm{}^{+i\mathrm{}}O_+^{(n)}\omega ^{n1}\frac{dn}{2\pi i},\hfill \\ & O_+^{(n)}=\left(i\frac{d}{dt}\right)^nO_+(t)|_{t=0}=_0^{\mathrm{}}O_+(\omega )\omega ^n𝑑\omega .\hfill \end{array}$$
(6.2)
As we shall see, not all of these relations survive renormalization.
We shall calculate matrix elements of these bare operators, therefore, we need the Feynman rules for them. If we retain $`i_+`$ in all brackets $`(iD_+)^n`$ in $`O_+^{(n)}`$, we obtain the quark–antiquark vertex with $`p_+^n`$ shown in Fig. 7. When transformed to the $`\omega `$-representation (6.2), this gives $`\delta (p_+\omega )`$, as expected from the form (6.1) of $`O_+(\omega )`$.
Now let’s retain a single $`A_+`$ in $`(iD_+)^n`$. After some combinatorics, this gives
$$(i_++A_+)^n\underset{m=1}{\overset{n}{}}\left(\begin{array}{c}n\\ m\end{array}\right)\left[(i_+)^{m1}A_+\right](i_+)^{nm}.$$
(6.3)
This gives the quark–antiquark–gluon vertex of $`O_+^{(n)}`$ with
$$\underset{m=1}{\overset{n}{}}\left(\begin{array}{c}n\\ m\end{array}\right)k_+^{m1}p_+^{nm}=\frac{(p_++k_+)^np_+^n}{k_+},$$
(6.4)
shown in Fig. 8. Transforming it to the $`\omega `$-representation (6.2) gives two $`\delta `$-functions (Fig. 8). Of course, the integral in $`d\omega `$ of this vertex vanishes, because $`O_+^{(0)}`$ does not interact with gluons.
The bare operators $`O_+(\omega )`$ can be expressed via the renormalized operators $`O_+(\omega ;\mu )`$ ($`\mu `$ is the $`\overline{\text{MS}}`$ renormalization scale), or vice versa:
$`O_+(\omega )`$ $`={\displaystyle Z_+(\omega ,\omega ^{};\mu )O_+(\omega ^{};\mu )𝑑\omega ^{}},`$ (6.5)
$`O_+(\omega ;\mu )`$ $`={\displaystyle Z_+^1(\omega ,\omega ^{};\mu )O_+(\omega ^{})𝑑\omega ^{}}.`$ (6.6)
The renormalization “matrices” $`Z_+(\omega ,\omega ^{};\mu )`$ and $`Z_+^1(\omega ,\omega ^{};\mu )`$ are inverse to each other:
$$\begin{array}{cc}\hfill Z_+(\omega ,\omega ^{\prime \prime };\mu )Z_+^1(\omega ^{\prime \prime },\omega ^{};\mu )𝑑\omega ^{\prime \prime }& =\delta (\omega \omega ^{}),\hfill \\ \hfill Z_+^1(\omega ,\omega ^{\prime \prime };\mu )Z_+(\omega ^{\prime \prime },\omega ^{};\mu )𝑑\omega ^{\prime \prime }& =\delta (\omega \omega ^{}).\hfill \end{array}$$
(6.7)
The renormalized operators $`O_+(\omega ;\mu )`$ obey the renormalization-group equation
$$\frac{O_+(\omega ;\mu )}{\mathrm{log}\mu }+\mathrm{\Gamma }_+(\omega ,\omega ^{};\mu )O_+(\omega ^{};\mu )𝑑\omega ^{}=0,$$
(6.8)
where the anomalous dimension “matrix” is
$$\begin{array}{cc}\hfill \mathrm{\Gamma }_+(\omega ,\omega ^{};\mu )& =Z_+^1(\omega ,\omega ^{\prime \prime };\mu )\frac{Z_+(\omega ^{\prime \prime },\omega ^{};\mu )}{\mathrm{log}\mu }𝑑\omega ^{\prime \prime }\hfill \\ & =\frac{Z_+^1(\omega ,\omega ^{\prime \prime };\mu )}{\mathrm{log}\mu }Z_+(\omega ^{\prime \prime },\omega ^{};\mu )𝑑\omega ^{\prime \prime }.\hfill \end{array}$$
(6.9)
This equation can be derived in two ways: either we differentiate (6.5) in $`d\mathrm{log}\mu `$ and obtain 0 (because the bare operator is $`\mu `$-independent), or we differentiate (6.6).
At one loop
$$\begin{array}{cc}\hfill Z_+(\omega ,\omega ^{};\mu )& =\delta (\omega \omega ^{})+z_+^{(1)}(\omega ,\omega ^{};\mu )a_s+\mathrm{}\hfill \\ \hfill Z_+^1(\omega ,\omega ^{};\mu )& =\delta (\omega \omega ^{})z_+^{(1)}(\omega ,\omega ^{};\mu )a_s+\mathrm{}\hfill \\ \hfill \mathrm{\Gamma }_+(\omega ,\omega ^{};\mu )& =\mathrm{\Gamma }_+^{(1)}(\omega ,\omega ^{};\mu )a_s+\mathrm{}\hfill \end{array}$$
(6.10)
where
$$a_s=\frac{\alpha _s(\mu )}{4\pi }.$$
From (6.9) we obtain
$$\mathrm{\Gamma }_+^{(1)}(\omega ,\omega ^{};\mu )=\frac{z_+^{(1)}(\omega ,\omega ^{};\mu )}{\mathrm{log}\mu }2\epsilon z_+^{(1)}(\omega ,\omega ^{};\mu ).$$
(6.11)
The matrix element of $`O_+(\omega )`$ between a state with the light quark with momentum $`p`$ and the heavy antiquark with momentum $`p^{}`$ (which are off-shell) and vacuum is
$$\begin{array}{cc}\hfill M=& <0|O_+(\omega )|q(p),Q_v^{}(p^{})>=Z_q^{1/2}\stackrel{~}{Z}_Q^{1/2}\left[\delta (p_+\omega )\gamma _++M_1+M_2+M_3\right]\hfill \\ \hfill =& <0|O_+(\omega ;\mu )|q(p),Q_v^{}(p^{})>+a_sz_+^{(1)}(\omega ,\omega ^{};\mu )\delta (p_+\omega ^{})\gamma _+𝑑\omega ^{},\hfill \end{array}$$
(6.12)
where $`M_{1,2,3}`$ are the contributions of the one-loop diagrams shown in Fig. 9, the matrix element of the renormalized operator is finite at $`\epsilon 0`$, and $`z_+^{(1)}(\omega ,\omega ^{};\mu )`$ contains only negative powers of $`\epsilon `$ (in the term with $`z_+^{(1)}`$, we may substitute $`\delta (p_+\omega ^{})\gamma _+`$ instead of $`<0|O_+(\omega ^{};\mu )|q(p),Q_v^{}(p^{})>`$). This allows us to find the renormalization “matrix” with one-loop accuracy.
Let’s calculate the first diagram (Fig. 10) for
$$p_+=\omega ^{},p_{}=0,p^2=p_+p_{}<0.$$
It is
$$M_1=iC_Fg_0^2\frac{d^dk}{(2\pi )^d}\frac{\delta (p_+\omega )\delta (k_+\omega )}{p_+k_+}\frac{\gamma _+\text{/}k\gamma _+}{\left[(pk)^2i0\right]\left[k^2i0\right]}.$$
(6.13)
The numerator can be simplified as
$$\gamma _+\text{/}k\gamma _+=2k_+\gamma _+.$$
This diagram can be written in the form
$$M_1=2C_F\frac{g_0^2(p^2)^\epsilon }{(4\pi )^{d/2}}\mathrm{\Gamma }(\epsilon )\gamma _+\left[f_1(\omega ,\omega ^{})\delta (\omega \omega ^{})f_1(\omega ^{\prime \prime },\omega ^{})𝑑\omega ^{\prime \prime }\right],$$
(6.14)
where the term with $`f_1(\omega ,\omega ^{})`$ comes from the second $`\delta `$-function in (6.13). This function is
$$\pi ^{d/2}(p^2)^\epsilon \mathrm{\Gamma }(\epsilon )f_1(\omega ,\omega ^{})=i\frac{\omega }{\omega ^{}\omega }d^dk\frac{\delta (k_+\omega )}{\left[(pk)^2i0\right]\left[k^2i0\right]}.$$
Now we use $`\alpha `$ parametrization
$$\frac{1}{k^2i0}=_0^{\mathrm{}}e^{(k^2+i0)\alpha }𝑑\alpha $$
(6.15)
for both denominators, and also write the $`\delta `$-function as
$$\delta (k_+\omega )=2\delta (2(k_+\omega ))=2_{\mathrm{}}^+\mathrm{}\mathrm{exp}\left[2(k_+\omega )\nu \right]\frac{d\nu }{2\pi }.$$
(6.16)
We obtain
$$\begin{array}{cc}& \pi ^{d/2}(p^2)^\epsilon \mathrm{\Gamma }(\epsilon )f_1(\omega ,\omega ^{})=2i\frac{\omega }{\omega ^{}\omega }d^dk𝑑\alpha _1𝑑\alpha _2\frac{d\nu }{2\pi }\hfill \\ & \times \mathrm{exp}\left[\alpha _1(kp)^2+\alpha _2k^2+2i\nu (kn_+\omega )\right].\hfill \end{array}$$
Shifting the integration momentum as
$$k^{}=k\frac{\alpha _1pi\nu n_+}{\alpha _1+\alpha _2}$$
to form a full square, we obtain
$$\begin{array}{cc}& 2i\frac{\omega }{\omega ^{}\omega }𝑑\alpha _1𝑑\alpha _2\mathrm{exp}\left[\frac{\alpha _1\alpha _2}{\alpha _1+\alpha _2}p^2\right]\hfill \\ & \times \frac{d\nu }{2\pi }\mathrm{exp}\left[2i\nu (\frac{\alpha _1}{\alpha _1+\alpha _2}\omega ^{}\omega )\right]d^dk^{}\mathrm{exp}\left[(\alpha _1+\alpha _2)(k^2+i0)\right]\hfill \end{array}$$
The integral in $`d^dk^{}`$ is calculated using the Wick rotation $`k_0=ik_{E0}`$:
$$d^dke^{\alpha (k^2+i0)}=id^dk_Ee^{\alpha k_E^2}=i\left(\frac{\pi }{\alpha }\right)^{d/2}.$$
(6.17)
The integral in $`d\nu `$ gives a $`\delta `$-function. We obtain an integral in two $`\alpha `$-parameters:
$$\begin{array}{cc}& (p^2)^\epsilon \mathrm{\Gamma }(\epsilon )f_1(\omega ,\omega ^{})\hfill \\ & =\frac{\omega }{\omega ^{}\omega }𝑑\alpha _1𝑑\alpha _2(\alpha _1+\alpha _2)^{d/2}\mathrm{exp}\left[\frac{\alpha _1\alpha _2}{\alpha _1+\alpha _2}p^2\right]\delta \left(\frac{\alpha _1}{\alpha _1+\alpha _2}\omega ^{}\omega \right).\hfill \end{array}$$
It is always possible to calculate one integral in a “radial” variable $`\eta `$ in the space of $`\alpha `$-parameters via $`\mathrm{\Gamma }`$-function. In this case, the most convenient choice of such a variable is $`\eta =\alpha _1+\alpha _2`$. Therefore, we insert $`\delta (\alpha _1+\alpha _2\eta )d\eta `$ under the integral sign, and substitute $`\alpha _i=\eta x_i`$:
$$\frac{\omega }{\omega ^{}\omega }𝑑x_1𝑑x_2\delta (x_1+x_21)\delta (x_1\omega ^{}\omega )𝑑\eta \eta ^{1+\epsilon }e^{(p^2)x_1x_2\eta }.$$
The final result is
$$f_1(\omega ,\omega ^{})=\frac{\theta (\omega ^{}\omega )}{(\omega ^{}\omega )^{1+\epsilon }}\frac{\omega ^{1\epsilon }}{(\omega ^{})^{12\epsilon }}.$$
(6.18)
Functions $`F(\omega ,\omega ^{})`$ which appear in the evolution kernel should be understood as distributions: they are always integrated with smooth test functions $`\phi (\omega ^{})`$. The distribution $`[F(\omega ,\omega ^{})]_+`$ is defined by
$$\left[F(\omega ,\omega ^{})\right]_+\phi (\omega ^{})𝑑\omega ^{}=F(\omega ,\omega ^{})\left(\phi (\omega ^{})\phi (\omega )\right)𝑑\omega ^{}.$$
(6.19)
Therefore, formally we can write
$$F(\omega ,\omega ^{})=\left[F(\omega ,\omega ^{})\right]_++\delta (\omega \omega ^{})F(\omega ,\omega ^{\prime \prime })𝑑\omega ^{\prime \prime }.$$
(6.20)
The result (6.14), (6.18) of the calculation of the diagram in Fig. 10 can be written via a $`+`$-distribution as
$$M_1=2C_F\frac{g_0^2(p^2)^\epsilon }{(4\pi )^{d/2}}\mathrm{\Gamma }(\epsilon )\gamma _+\left[\left[f_1(\omega ,\omega ^{})\right]_++\delta (\omega \omega ^{})\left(f_1(\omega ,\omega ^{\prime \prime })𝑑\omega ^{\prime \prime }f_1(\omega ^{\prime \prime },\omega )𝑑\omega ^{\prime \prime }\right)\right].$$
(6.21)
The coefficient of $`\delta (\omega \omega ^{})`$ here can be calculated by substitution $`x=\omega ^{\prime \prime }/\omega `$:
$$_1^{\mathrm{}}x^{1+2\epsilon }(1x)^{1\epsilon }𝑑x_0^1x^{1\epsilon }(1x)^{1\epsilon }𝑑x.$$
Substituting $`x1/x`$ in the first integral, we have
$$_0^1\left(x^\epsilon x^{1\epsilon }\right)(1x)^{1\epsilon }𝑑x=_0^1x^\epsilon (1x)^\epsilon 𝑑x.$$
Therefore, the final result for $`M_1`$ (Fig. 10) is
$$M_1=2C_F\frac{g_0^2(p^2)^\epsilon }{(4\pi )^{d/2}}\mathrm{\Gamma }(\epsilon )\gamma _+\left[\left(\frac{\theta (\omega ^{}\omega )}{(\omega ^{}\omega )^{1+\epsilon }}\frac{\omega ^{1\epsilon }}{(\omega ^{})^{12\epsilon }}\right)_++\frac{\mathrm{\Gamma }^2(1\epsilon )}{\mathrm{\Gamma }(22\epsilon )}\delta (\omega \omega ^{})\right].$$
(6.22)
Now we shall calculate the second diagram (Fig. 11) for
$$p^{}v=\omega _1<0.$$
It is
$$M_2=iC_Fg_0^2\frac{d^dk}{(2\pi )^d}\frac{\delta (p_++k_+\omega )\delta (p_+\omega )}{k_+}\frac{v_+\gamma _+}{\left[k^2i0\right]\left[(p^{}k)vi0\right]}.$$
(6.23)
This diagram can be written in the form
$$M_2=2C_F\frac{g_0^2}{(4\pi )^{d/2}}\mathrm{\Gamma }(\epsilon )\gamma _+\left[f_2(\omega \omega ^{})\delta (\omega \omega ^{})f_2(\omega ^{\prime \prime })𝑑\omega ^{\prime \prime }\right],$$
(6.24)
where the term with $`f_2(\omega \omega ^{})`$ comes from the first $`\delta `$-function in (6.23). This function is
$$\pi ^{d/2}\mathrm{\Gamma }(\epsilon )f_2(\omega ^{\prime \prime })=\frac{i}{2\omega ^{\prime \prime }}d^dk\frac{\delta (k_+\omega ^{\prime \prime })}{\left[k^2i0\right]\left[(p^{}k)vi0\right]}.$$
Now we use $`\alpha `$-parametrization (6.15) for both denominators (it is convenient to multiply the linear denominator by 2 before this) and (6.16) for the $`\delta `$-function, and obtain
$$\frac{2i}{\omega ^{\prime \prime }}d^dk𝑑\alpha _1𝑑\alpha _2\frac{d\nu }{2\pi }\mathrm{exp}\left[\alpha _1k^2+2\alpha _2(p^{}k)v+2i\nu (kn_+\omega ^{\prime \prime })\right].$$
Shifting the integration momentum as
$$k^{}=k\frac{\alpha _2vi\nu n_+}{\alpha _1}$$
to form a full square, we obtain
$$\begin{array}{cc}& \frac{2i}{\omega ^{\prime \prime }}𝑑\alpha _1𝑑\alpha _2\mathrm{exp}\left[\frac{\alpha _2^2}{\alpha _1}+2\omega _1\alpha _2\right]\hfill \\ & \times \frac{d\nu }{2\pi }\mathrm{exp}\left[2i\nu (\frac{\alpha _2}{\alpha _1}\omega ^{\prime \prime })\right]d^dk^{}\mathrm{exp}\left[\alpha _1(k^2+i0)\right].\hfill \end{array}$$
Taking the integrals in $`d^dk^{}`$ (6.17) and $`d\nu `$, we obtain the integral in two $`\alpha `$-parameters:
$$\mathrm{\Gamma }(\epsilon )f_2(\omega ^{\prime \prime })=\frac{1}{\omega ^{\prime \prime }}𝑑\alpha _1𝑑\alpha _2\alpha _1^{d/2}\mathrm{exp}\left[\frac{\alpha _2^2}{\alpha _1}+2\omega _1\alpha _2\right]\delta \left(\frac{\alpha _2}{\alpha _1}\omega ^{\prime \prime }\right).$$
Now the best choice of the “radial” variable is $`\alpha _1`$, so we simply substitute $`\alpha _2=\alpha _1y`$:
$$\frac{1}{\omega ^{\prime \prime }}𝑑y\delta (y\omega ^{\prime \prime })𝑑\alpha _1\alpha _1^{1+\epsilon }e^{y(y2\omega _1)\alpha _1}.$$
Finally,
$$f_2(\omega ^{\prime \prime })=\frac{\theta (\omega ^{\prime \prime })}{(\omega ^{\prime \prime })^{1+\epsilon }(\omega ^{\prime \prime }2\omega _1)^\epsilon }.$$
(6.25)
Now we rewrite $`M_2`$ (6.24), (6.25) via a $`+`$-distribution:
$$M_2=2C_F\frac{g_0^2}{(4\pi )^{d/2}}\mathrm{\Gamma }(\epsilon )\gamma _+\left[\left[f_2(\omega \omega ^{})\right]_++\delta (\omega \omega ^{})\left(_0^\omega f_2(\omega \omega ^{\prime \prime })𝑑\omega ^{\prime \prime }_0^{\mathrm{}}f_2(\omega ^{\prime \prime })𝑑\omega ^{\prime \prime }\right)\right].$$
(6.26)
The coefficient of $`\delta (\omega \omega ^{})`$ is
$$_\omega ^{\mathrm{}}f_2(\omega ^{\prime \prime })𝑑\omega ^{\prime \prime }.$$
(6.27)
We introduced $`\omega _1`$ only to regularize possible infrared problems; we only need the UV divergence of this diagram, which does not depend on $`\omega _1`$. Therefore, we may assume $`|\omega _1|\omega `$. The coefficient of $`\delta (\omega \omega ^{})`$ is
$$_\omega ^{\mathrm{}}(\omega ^{\prime \prime })^{12\epsilon }𝑑\omega ^{\prime \prime }=\frac{\omega ^{2\epsilon }}{2\epsilon },$$
(6.28)
and the final result for $`M_2`$ (Fig. 11) is
$$M_2=2C_F\frac{g_0^2}{(4\pi )^{d/2}}\mathrm{\Gamma }(\epsilon )\gamma _+\left[\left(\frac{\theta (\omega \omega ^{})}{(\omega \omega ^{})^{1+2\epsilon }}\right)_+\frac{\omega ^{2\epsilon }}{2\epsilon }\delta (\omega \omega ^{})\right].$$
(6.29)
This one-loop diagram contains a $`1/\epsilon ^2`$ UV divergence! It is in the coefficient of $`\delta (\omega \omega ^{})`$ given by the integral (6.27). The function $`f_2(\omega ^{\prime \prime })`$ has a $`1/\epsilon `$ UV divergence in the integration in transverse momentum; the longitudinal integral (6.27) is again UV divergent (6.28). The operator at the vertex of Fig. 11 is (after integration in $`\omega `$) a light-like Wilson line; the heavy-antiquark propagator is a time-like Wilson line. The vertex correction to a time-like – light-like cusp on a Wilson line is known to have a $`1/\epsilon ^2`$ UV divergence at one loop .
The third diagram (Fig. 12) is
$$M_3=iC_Fg_0^2\frac{d^dk}{(2\pi )^d}\frac{\delta (k_+\omega )\gamma _+\text{/}k\text{/}v}{\left[(kp)^2i0\right]\left[k^2i0\right]\left[(p^{}+pk)vi0\right]}.$$
(6.30)
The numerator can be simplified as
$$\gamma _+\text{/}k\text{/}v=k_+\gamma _+.$$
We shall now demonstrate that this diagram is UV-finite and hence does not contribute to the renormalization constant.
Using $`\alpha `$-parametrization (6.15) for the denominators and (6.16) for the $`\delta `$-function, and obtain
$$\begin{array}{cc}\hfill M_3=& 4iC_Fg_0^2\omega \gamma _+d^dk𝑑\alpha _1𝑑\alpha _2𝑑\alpha _3\frac{d\nu }{2\pi }\hfill \\ & \times \mathrm{exp}\left[\alpha _1(kp)^2+\alpha _2k^2+2\alpha _3(p+p^{}k)v+2i\nu (kn_+\omega )\right].\hfill \end{array}$$
Shifting the integration momentum as
$$k^{}=k\frac{\alpha _1p+\alpha _3vi\nu n_+}{\alpha _1+\alpha _2},$$
we obtain
$$\begin{array}{cc}\hfill M_3=& 4iC_Fg_0^2\omega \gamma _+𝑑\alpha _1𝑑\alpha _2𝑑\alpha _3e^A\hfill \\ & \times \frac{d\nu }{2\pi }\mathrm{exp}\left[2i\nu (\frac{\alpha _1\omega ^{}+\alpha _3}{\alpha _1+\alpha _2}\omega )\right]\frac{d^dk^{}}{(2\pi )^d}e^{(\alpha _1+\alpha _2)(k^2+i0)}\hfill \\ \hfill =& 2C_F\frac{g_0^2}{(4\pi )^{d/2}}\omega \gamma _+𝑑\alpha _1𝑑\alpha _2𝑑\alpha _3\delta \left(\frac{\alpha _1\omega ^{}+\alpha _3}{\alpha _1+\alpha _2}\omega \right)e^A,\hfill \end{array}$$
where
$$A=\frac{\alpha _3(\alpha _3\alpha _2\omega ^{})+\alpha _2(\alpha _1\omega ^{}+\alpha _3)(p^2)}{\alpha _1+\alpha _2}+\alpha _3(2\omega ).$$
Inserting $`\delta (\alpha _1+\alpha _2\eta )d\eta `$ under the integral sign and making substitutions $`\alpha _{1,2}=\eta x_{1,2}`$, $`\alpha _3=\eta y`$, we have $`A=a\eta `$, and the integral in the “radial variable” $`\eta `$ is easily calculated:
$$\begin{array}{cc}\hfill M_3& =2C_F\frac{g_0^2}{(4\pi )^{d/2}}\omega \gamma _+𝑑x_1𝑑y\delta (y+\omega ^{}x_1\omega )𝑑\eta \eta ^\epsilon e^{a\eta }\hfill \\ & =2C_F\frac{g_0^2}{(4\pi )^{d/2}}\mathrm{\Gamma }(1+\epsilon )\frac{\omega }{\omega ^{}}\gamma _+_{\mathrm{max}(0,\omega \omega ^{})}^\omega \frac{dy}{a^{1+\epsilon }},\hfill \end{array}$$
where
$$a=\left[\omega \omega ^{}+\frac{\omega }{\omega ^2}(p^2)2\omega _1\right]y+(\omega \omega ^{})\frac{\omega }{\omega ^2}p^2.$$
The integral in $`y`$ is finite at $`\epsilon 0`$ (and easy to calculate).
Now we are ready to find the renormalization constant $`Z_+(\omega ,\omega ^{};\mu )`$. Re-expressing
$$\frac{g_0^2}{(4\pi )^{d/2}}=a_s\mu ^{2\epsilon }e^{\gamma _E\epsilon }$$
in the matrix element (6.12), we have
$$\begin{array}{cc}\hfill M=& \gamma _+\left(Z_q\stackrel{~}{Z}_Q\right)^{1/2}[\delta (\omega \omega ^{})+2C_F\frac{\alpha _s(\mu )}{4\pi \epsilon }\hfill \\ & \times (\left(\frac{\theta (\omega ^{}\omega )}{\omega ^{}\omega }\frac{\omega }{\omega ^{}}\right)_++\delta (\omega \omega ^{})\hfill \\ & +\left(\frac{\theta (\omega \omega ^{})}{\omega \omega ^{}}\right)_+(\frac{1}{2\epsilon }\mathrm{log}\frac{\omega }{\mu })\delta (\omega \omega ^{})+𝒪(\epsilon ))].\hfill \end{array}$$
(6.31)
The quark-field renormalization constants in QCD and HQET th the Feynman gauge are
$$Z_q=1C_F\frac{\alpha _s}{4\pi \epsilon },\stackrel{~}{Z}_Q=1+2C_F\frac{\alpha _s}{4\pi \epsilon }.$$
(6.32)
Finally, the renormalization constant “matrix” at one loop is
$$\begin{array}{cc}& Z_+(\omega ,\omega ^{};\mu )=\delta (\omega \omega ^{})+2C_F\frac{\alpha _s(\mu )}{4\pi \epsilon }\hfill \\ & \times \left[\left(\frac{\theta (\omega ^{}\omega )}{\omega ^{}\omega }\frac{\omega }{\omega ^{}}+\frac{\theta (\omega \omega ^{})}{\omega \omega ^{}}\right)_++\left(\frac{1}{2\epsilon }+\mathrm{log}\frac{\omega }{\mu }+\frac{5}{4}\right)\delta (\omega \omega ^{})\right].\hfill \end{array}$$
(6.33)
The evolution kernel (6.9) has the structure
$$\mathrm{\Gamma }_+(\omega ,\omega ^{};\mu )=\mathrm{\Gamma }(\omega ,\omega ^{};a_s)+\left[\mathrm{\Gamma }(\omega ,\omega ^{};a_s)\mathrm{log}\frac{\omega }{\mu }+\stackrel{~}{\gamma }_j(a_s)+\gamma (a_s)\right].\delta (\omega \omega ^{}).$$
(6.34)
The logarithm $`\mathrm{log}(\omega /\mu )`$ only appears linearly, to all orders of perturbation theory; the coefficient of this logarithm is the cusp anomalous dimension :
$$\mathrm{\Gamma }(a_s)=\mathrm{\Gamma }_0a_s+\mathrm{\Gamma }_1a_s^2+\mathrm{}\mathrm{\Gamma }_0=4C_F$$
(6.35)
(the two-loop term is also known ). The non-logarithmic part is written as $`\stackrel{~}{\gamma }_j+\gamma `$, where the anomalous dimension of the local current $`j`$ (3.1) is
$$\stackrel{~}{\gamma }_j(a_s)=\stackrel{~}{\gamma }_{j0}a_s+\stackrel{~}{\gamma }_{j1}a_s^2+\mathrm{}\stackrel{~}{\gamma }_{j0}=3C_F$$
(6.36)
(the two- and three-loop terms are also known ). It determines the evolution of $`F(\mu )`$; the difference
$$\gamma (a_s)=\gamma _0a_s+\gamma _1a_s^2+\mathrm{}\gamma _0=2C_F$$
(6.37)
appears in the evolution equation for $`\phi _+(\omega ;\mu )`$. The non-$`\delta `$ term is
$$\mathrm{\Gamma }(\omega ,\omega ^{},a_s)=\mathrm{\Gamma }_0(\omega ,\omega ^{})a_s+\mathrm{}\mathrm{\Gamma }_0(\omega ,\omega ^{})=\left(\frac{\theta (\omega ^{}\omega )}{\omega ^{}\omega }\frac{\omega }{\omega ^{}}+\frac{\theta (\omega \omega ^{})}{\omega \omega ^{}}\right)_+.$$
(6.38)
The evolution kernel has dimensionality of $`1/\text{energy}`$.
The evolution equation for the distribution amplitude is
$$\frac{\phi _+(\omega ;\mu )}{\mathrm{log}\mu }+\left[\mathrm{\Gamma }(a_s)\mathrm{log}\frac{\omega }{\mu }+\gamma (a_s)\right]\phi _+(\omega ;\mu )+\mathrm{\Gamma }(\omega ,\omega ^{};a_s)\phi _+(\omega ^{};\mu )𝑑\omega ^{}=0.$$
(6.39)
How to solve it? Powers $`\omega ^n`$ are eigenfunctions of the integral operator in (6.39), by dimensionality:
$$\mathrm{\Gamma }(\omega ,\omega ^{};a_s)\omega ^n𝑑\omega ^{}=\stackrel{~}{\mathrm{\Gamma }}(n,a_s)\omega ^n,$$
(6.40)
where
$$\stackrel{~}{\mathrm{\Gamma }}(n,a_s)=\stackrel{~}{\mathrm{\Gamma }}_0a_s+\mathrm{}\stackrel{~}{\mathrm{\Gamma }}_0(n)=4C_F\left[\psi (1+n)+\psi (1n)+2\gamma _E\right],$$
(6.41)
from (6.38) (here $`\gamma _E`$ is the Euler constant). They are not eigenfunctions of the whole evolution operator, due to the logarithmic term. We can construct solutions with the power depending on $`\mu `$:
$$\frac{}{\mathrm{log}\mu }\left(\frac{\omega }{\mu _0}\right)^{n+\xi (\mu )}=\left(\frac{\omega }{\mu _0}\right)^{n+\xi (\mu )}\frac{d\xi }{d\mathrm{log}\mu }\mathrm{log}\frac{\omega }{\mu _0}.$$
If the function $`\xi (\mu )`$ obeys
$$\frac{d\xi }{d\mathrm{log}\mu }=\mathrm{\Gamma }(a_s),$$
(6.42)
then $`\mathrm{log}\omega `$ will cancel in the evolution equation.
Dividing the definition (6.42) by
$$\frac{d\mathrm{log}a_s}{d\mathrm{log}\mu }=2\beta (a_s),\beta (a_s)=\beta _0a_s+\beta _1a_s^2+\mathrm{}$$
(6.43)
and integrating, we obtain
$$\xi =_{a_{s0}}^{a_s}\frac{\mathrm{\Gamma }(a_s)}{2\beta (a_s)}\frac{da_s}{a_s}=\frac{\mathrm{\Gamma }_0}{2\beta _0}\left[\mathrm{log}\frac{a_s}{a_{s0}}+\left(\frac{\mathrm{\Gamma }_1}{\mathrm{\Gamma }_0}\frac{\beta _1}{\beta _0}\right)\left(a_sa_{s0}\right)+\mathrm{}\right],$$
(6.44)
where $`a_{s0}=\alpha _s(\mu _0)/(4\pi )`$. We shall use $`\xi `$ as an independent variable (“time”) in the evolution equation instead of $`\mu `$. The function
$$\phi (\omega ,\xi )=\left(\frac{\omega }{\mu _0}\right)^{n+\xi }e^{U(\xi )}$$
(6.45)
satisfies the evolution equation
$$\frac{\mathrm{log}\phi }{\xi }\mathrm{log}\frac{\omega }{\mu _0}+\frac{\stackrel{~}{\mathrm{\Gamma }}(n+\xi ,a_s)+\gamma (a_s)}{\mathrm{\Gamma }(a_s)}=0$$
(6.46)
if
$$\frac{dU}{d\xi }=\mathrm{log}\frac{\mu }{\mu _0}\frac{\stackrel{~}{\mathrm{\Gamma }}(n+\xi ,a_s)+\gamma (a_s)}{\mathrm{\Gamma }(a_s)}.$$
At the leading order,
$$a_s=a_{s0}\mathrm{exp}\left(\frac{2\beta _0}{\mathrm{\Gamma }_0}\xi \right),\mathrm{log}\frac{\mu }{\mu _0}=\frac{2\pi }{\beta _0}\left(\frac{1}{a_s}\frac{1}{a_{s0}}\right)=\frac{2\pi }{\beta _0a_{s0}}\left[\mathrm{exp}\left(\frac{2\beta _0}{\mathrm{\Gamma }_0}\xi \right)1\right],$$
and
$$U(\xi )=\frac{2\pi }{\beta _0a_{s0}}\left[\frac{\mathrm{\Gamma }_0}{2\beta _0}\left(\mathrm{exp}\left(\frac{2\beta _0}{\mathrm{\Gamma }_0}\xi \right)1\right)\xi \right]_0^\xi \frac{\stackrel{~}{\mathrm{\Gamma }_0}(n+\xi )}{\mathrm{\Gamma }_0}𝑑\xi \frac{\gamma _0}{\mathrm{\Gamma }_0}\xi .$$
Here, from (6.41),
$$_0^\xi \frac{\stackrel{~}{\mathrm{\Gamma }_0}(n+\xi )}{\mathrm{\Gamma }_0}𝑑\xi =\mathrm{log}\frac{\mathrm{\Gamma }(1+n+\xi )\mathrm{\Gamma }(1n)}{\mathrm{\Gamma }(1n\xi )\mathrm{\Gamma }(1+n)}+2\gamma _E\xi .$$
Finally, at the leading order,
$$e^{U(\xi )}=\left(\frac{\mathrm{\Lambda }_{\overline{\text{MS}}}}{\mu _0}\right)^{\frac{\mathrm{\Gamma }_0}{2\beta _0}\left[\mathrm{exp}\left(\frac{2\beta _0}{\mathrm{\Gamma }_0}\xi \right)1\right]\xi }\frac{\mathrm{\Gamma }(1n\xi )\mathrm{\Gamma }(1+n)}{\mathrm{\Gamma }(1+n+\xi )\mathrm{\Gamma }(1n)}\mathrm{exp}\left[\left(\frac{\gamma _0}{\mathrm{\Gamma }_0}+2\gamma _E\right)\xi \right].$$
(6.47)
The distribution amplitude at $`\mu _0`$ can be expressed via its moments (4.20):
$$\phi _+(\omega ;\mu _0)=_i\mathrm{}^{+i\mathrm{}}<\omega ^{1n}>_+^{(\mu _0)}\omega ^n\frac{dn}{2\pi i}.$$
Substituting the solutions (6.45) instead of $`\omega ^n`$ gives us the solution of the evolution equation (6.39) with initial conditions. At the leading order (6.47),
$$\begin{array}{cc}\hfill \phi _+(\omega ;\mu )=& \left(\frac{\mathrm{\Lambda }_{\overline{\text{MS}}}}{\mu _0}\right)^{\frac{\mathrm{\Gamma }_0}{2\beta _0}\left[\mathrm{exp}\left(\frac{2\beta _0}{\mathrm{\Gamma }_0}\xi \right)1\right]}\mathrm{exp}\left[\left(\frac{\gamma _0}{\mathrm{\Gamma }_0}+2\gamma _E\right)\xi \right]\left(\frac{\omega }{\mathrm{\Lambda }_{\overline{\text{MS}}}}\right)^\xi \hfill \\ & \times _{\mathrm{}}^+\mathrm{}<\omega ^{1in}>_+^{(\mu _0)}\omega ^{in}\frac{\mathrm{\Gamma }(1n\xi )\mathrm{\Gamma }(1+n)}{\mathrm{\Gamma }(1+n+\xi )\mathrm{\Gamma }(1n)}\frac{dn}{2\pi }.\hfill \end{array}$$
(6.48)
Qualitatively, each “bin” in the distribution amplitude at $`\mu _0`$ near some finite $`\omega ^{}`$ produces, after evolution to a larger $`\mu `$, a radiative tail, slowly decreasing (as $`1/\omega `$) at $`\omega \omega ^{}`$, due to the evolution kernel $`\mathrm{\Gamma }(\omega ,\omega ^{};a_s)`$ (which behaves as $`1/\omega `$ at large $`\omega `$, by dimensionality). Therefore, the behaviour of the distribution amplitude at large $`\omega `$ is $`1/\omega `$, up to logarithms. The normalization integral (4.14) of the distribution amplitude logarithmically diverges at large $`\omega `$; its moments (4.18) with $`n>0`$ are power-divergent. This means that the expression (4.17) for the local operators $`O_+^{(n)}`$ ($`n0`$) via $`O_+(\omega )`$ is not valid for the renormalized operators. Renormalization of $`O_+(\omega )`$ removes UV divergences in transverse momenta; in order to renormalize $`O_+^{(n)}`$, we should also remove longitudinal UV divergences (at large $`\omega `$).
As an illustration , let’s suppose that at a low $`\mu _0`$ the distribution amplitude $`\phi _+(\omega )`$ is given by the simple model (7.24). Namely, this initial condition is taken at the scale $`\mu _0`$ where
$$\alpha _s(\mu _0)=1\frac{\mu _0}{\mathrm{\Lambda }_{\overline{\text{MS}}}}=\mathrm{exp}\frac{2\pi }{\beta _0},$$
where the quark model is supposed to work, so that the light quark has only momenta of order $`\overline{\mathrm{\Lambda }}`$ (no radiative tail). This function is shown by the solid line in Fig. (13), where $`\omega `$ is measured in units of $`\omega _0`$ (7.25). Using the leading-order solution (6.48) of the evolution equation (6.39) and supposing $`\overline{\mathrm{\Lambda }}/\mathrm{\Lambda }_{\overline{\text{MS}}}=1.25`$, we obtain the distribution amplitude at the scale where $`\alpha _s(\mu )=0.5`$ (dashed line) and $`0.3`$ (dashed-dotted line). We can see that the main part of the distribution amplitude (at $`\omega \overline{\mathrm{\Lambda }}`$) becomes lower and the radiative tail becomes more prominent when $`\mu `$ grows.
## 7 Sum rules
In order to estimate the quark–antiquark distribution amplitudes $`\phi _\pm (\omega )`$ of $`B`$-meson from QCD sum rules, we consider the correlator of the local current $`j_+`$ (3.2) (having the quantum numbers of the ground-state meson) and the bilocal operator $`\stackrel{~}{O}_\pm (t)`$ (4.11):
$$i<T\stackrel{~}{O}_\pm (t)\overline{ȷ}_+(x)>=\gamma _\pm \frac{1+\gamma ^0}{2}\delta (\stackrel{}{x})\theta (x^0)\stackrel{~}{\mathrm{\Pi }}_\pm (x^0,t).$$
(7.1)
We are most interested in its Fourier transform
$$\mathrm{\Pi }_\pm (x^0,\omega )=\stackrel{~}{\mathrm{\Pi }}_\pm (x^0,t)e^{i\omega t}\frac{dt}{2\pi }.$$
(7.2)
Analytically continuing it from $`x^0>0`$ to $`x^0=i\tau `$, we obtain
$$\mathrm{\Pi }_\pm (\tau ,\omega )=\rho _\pm (\epsilon ,\omega )e^{\epsilon \tau }𝑑\epsilon =F^2\stackrel{~}{\phi }_\pm (\omega )e^{\overline{\mathrm{\Lambda }}\tau }+\mathrm{\Pi }_\pm ^c(\tau ,\omega ),$$
(7.3)
where $`\rho _\pm (\epsilon ,\omega )`$ is the spectral density. The correlator contains the contribution of the ground-state meson (written explicitly in (7.3)) and of the continuum of excited states.
For sufficiently small $`\tau `$, we can calculate this correlator theoretically. The perturbative contribution (Fig. 14) is given in the fixed-point gauge $`x^\mu A_\mu (x)=0`$ by the light-quark propagator from $`x`$ to $`z`$:
$$\stackrel{~}{\mathrm{\Pi }}_+^{(1)}(\tau ,t)=\frac{N_c}{2\pi ^2\tau (\tau +2it)^2},\stackrel{~}{\mathrm{\Pi }}_{}^{(1)}(\tau ,t)=\frac{N_c}{2\pi ^2\tau ^2(\tau +2it)}.$$
(7.4)
Its Fourier transform is
$$\mathrm{\Pi }_+^{(1)}(\tau ,\omega )=\frac{N_c}{8\pi ^2\tau }\omega e^{\omega \tau /2},\mathrm{\Pi }_{}^{(1)}(\tau ,\omega )=\frac{N_c}{4\pi ^2\tau ^2}e^{\omega \tau /2}.$$
(7.5)
Inverting the Laplace transform (7.3), we find the spectral densities (the integration contour should be to the right of the singularity at $`\tau =0`$):
$$\rho _\pm (\epsilon ,\omega )=_i\mathrm{}^{+i\mathrm{}}\mathrm{\Pi }_\pm (\tau ,\omega )e^{\epsilon \tau }\frac{d\tau }{2\pi i}.$$
(7.6)
We obtain
$$\rho _+^{(1)}(\epsilon ,\omega )=\frac{N_c}{8\pi ^2}\omega \theta \left(\epsilon \frac{\omega }{2}\right),\rho _{}^{(1)}(\epsilon ,\omega )=\frac{N_c}{4\pi ^2}\left(\epsilon \frac{\omega }{2}\right)\theta \left(\epsilon \frac{\omega }{2}\right).$$
(7.7)
The perturbative spectral density $`\rho _\pm ^{(1)}(\epsilon ,\omega )`$ describes the contribution of the on-shell quark–antiquark intermediate state with energy $`\epsilon `$ into the correlator (7.3). The on-shell heavy antiquark has zero energy; the maximum value of $`p_+`$ of the on-shell light quark with energy $`\epsilon `$ is $`2\epsilon `$. This explains the $`\theta `$-functions in (7.7).
The quark condensate contribution is also important here. It contains the vacuum average
$$<\overline{q}(x)[x,0][0,z]q(z)>.$$
This non-collinear quark condensate can be expanded in terms of bilocal condensates , the leading term in this expansion is the bilocal quark condensate
$$<\overline{q}(0)[0,x]q(x)>=<\overline{q}q>f_S(x^2)$$
(7.8)
with $`xx+z`$. This leading term produces the same contribution into $`\stackrel{~}{\mathrm{\Pi }}_\pm `$:
$$\stackrel{~}{\mathrm{\Pi }}_\pm ^{(2)}(x^0,t)=\frac{1}{4}<\overline{q}q>f_S((x+z)^2).$$
(7.9)
The expansion of the bilocal quark condensate at small $`x`$ is
$$f_S(x^2)=1+\frac{m_0^2}{16}x^2+\mathrm{}$$
(7.10)
where
$$<\overline{q}G_{\mu \nu }\sigma ^{\mu \nu }q>=m_0^2<\overline{q}q>.$$
(7.11)
In general, it can be written as
$$f_S(x^2)=\stackrel{~}{f}_S(\nu )e^{\nu x^2}𝑑\nu ,$$
(7.12)
where $`\stackrel{~}{f}_S(\nu )`$ has the meaning of the virtuality distribution function of vacuum quarks. Its moments are expressed via vacuum averages of local operators:
$$\stackrel{~}{f}_S(\nu )𝑑\nu =1,\stackrel{~}{f}_S(\nu )\nu 𝑑\nu =\frac{m_0^2}{16},\mathrm{}$$
(7.13)
In terms of this function, the quark-condensate contribution into the Fourier-transformed correlators is
$$\mathrm{\Pi }_\pm ^{(2)}(\tau ,\omega )=\frac{<\overline{q}q>}{8\tau }\stackrel{~}{f}_S\left(\frac{\omega }{2\tau }\right)e^{\omega \tau /2}.$$
(7.14)
In other words, the virtuality distribution function appears directly in the sum rules for the $`B`$-meson distribution amplitudes! This is similar to the case of non-diagonal sum rules for the pion distribution amplitude .
The local operator expansion (7.10) gives
$$\stackrel{~}{f}_S(\nu )=\delta (\nu )\frac{m_0^2}{16}\delta ^{}(\nu )+\mathrm{}$$
(7.15)
Of course, a smooth function can always be expanded in derivatives of $`\delta `$-function, but such an expansion does not tell us much about the shape of this function, unless we sum an infinite number of terms. The bilocal quark condensate (7.8) at large $`x^2`$ behaves as
$$f_S(x^2)e^{\overline{\mathrm{\Lambda }}\sqrt{x^2}},$$
(7.16)
because the Wilson line $`[0,x]`$ can be considered an HQET heavy-quark propagator, and this vacuum average is the correlator of two HQET heavy–light currents, having $`B`$-meson as the lowest-energy intermediate state. Of course, the conditions (7.13), (7.16) don’t determine the shape of the distribution function $`\stackrel{~}{f}_S(\nu )`$ in a unique way. They are satisfied by the Bakulev–Mikhailov ansatz
$$\stackrel{~}{f}_S(\nu )=N\mathrm{exp}\left(\frac{\overline{\mathrm{\Lambda }}^2}{4\nu }\sigma \nu \right),$$
(7.17)
where the parameters $`\sigma `$, $`N`$ are determined by (7.13).
Now we equate the theoretical result for the correlators to the result obtained from a phenomenological model of the spectral densities. They have the contribution $`\delta (\epsilon \overline{\mathrm{\Lambda }})`$ of the ground-state meson and that of the continuum of excited states. As usual, this contribution is modeled by the perturbative spectral densities (7.7) starting from a continuum threshold energy $`\epsilon _c`$:
$$\rho _\pm (\epsilon ,\omega )=F^2\phi _\pm (\omega )\delta \left(\epsilon \overline{\mathrm{\Lambda }}\right)+\rho _\pm ^{(1)}(\epsilon ,\omega )\theta (\epsilon \epsilon _c).$$
(7.18)
With this model, the equality of the phenomenological expression for the correlator and the theoretical one,
$$F^2\phi _\pm (\omega )e^{\overline{\mathrm{\Lambda }}\tau }+_{\epsilon _c}^{\mathrm{}}\rho _\pm ^{(1)}(\epsilon ,\omega )e^{\epsilon \tau }𝑑\epsilon =_0^{\mathrm{}}\rho _\pm ^{(1)}(\epsilon ,\omega )e^{\epsilon \tau }𝑑\epsilon +\mathrm{\Pi }_\pm ^{(2)}(\tau ,\omega ),$$
becomes
$$F^2\phi _\pm (\omega )e^{\overline{\mathrm{\Lambda }}\tau }=_0^{\epsilon _c}\rho _\pm ^{(1)}(\epsilon ,\omega )e^{\epsilon \tau }𝑑\epsilon +\mathrm{\Pi }_\pm ^{(2)}(\tau ,\omega ),$$
(7.19)
where the perturbative spectral density is integrated over the “duality interval” of the ground-state meson (from $`0`$ to the continuum threshold $`\epsilon _c`$). We obtain the sum rules for the $`B`$-meson distribution amplitudes
$$\begin{array}{cc}& F^2\phi _+(\omega )e^{\overline{\mathrm{\Lambda }}\tau }=\frac{N_c}{8\pi ^2\tau }\omega e^{\omega \tau /2}\delta _0\left(\left(\epsilon _c\frac{\omega }{2}\right)\tau \right)\frac{<\overline{q}q>}{8\tau }\stackrel{~}{f}_S\left(\frac{\omega }{2\tau }\right)e^{\omega \tau /2},\hfill \\ & F^2\phi _{}(\omega )e^{\overline{\mathrm{\Lambda }}\tau }=\frac{N_c}{4\pi ^2\tau ^2}e^{\omega \tau /2}\delta _1\left(\left(\epsilon _c\frac{\omega }{2}\right)\tau \right)\frac{<\overline{q}q>}{8\tau }\stackrel{~}{f}_S\left(\frac{\omega }{2\tau }\right)e^{\omega \tau /2}.\hfill \end{array}$$
(7.20)
where the functions
$$\delta _n(x)=\theta (x)\left(1e^x\underset{m=0}{\overset{n}{}}\frac{x^m}{m!}\right)$$
(7.21)
describe the effect of subtracting continuum from the perturbative contribution when the spectral density is $`\epsilon ^n`$. The perturbative contributions vanish at $`\omega >2\epsilon _c`$, because of the properties of the spectral densities (7.7).
Setting $`t=0`$ in the correlator (7.1), or integrating (7.2) in $`d\omega `$, we obtain the well-known sum rule for $`F^2`$:
$$F^2e^{\overline{\mathrm{\Lambda }}\tau }=\frac{N_c}{2\pi ^2\tau ^3}\delta _2(\epsilon _c\tau )\frac{1}{4}<\overline{q}q>f_S(\tau ^2).$$
(7.22)
The natural energy scale in this sum rule is
$$k=\left(\frac{\pi ^2}{2N_c}<\overline{q}q\right)^{1/3}260\text{MeV}.$$
(7.23)
With $`m_0/(4k)=0.85`$, one finds the optimal value of the continuum threshold $`\epsilon _c/k=3`$; then a wide plato at $`1/(k\tau )[1.7,2.5]`$ exists, and yields $`\overline{\mathrm{\Lambda }}/k=1.65`$. We divide the sum rules (7.20) by (7.22) and take $`1/(k\tau )=2`$ (in the middle of the plato). The results are shown in Fig. 15, 16, where $`\omega `$ is measured in units of $`k`$. They are automatically normalized (4.14). Of course, the leading-order sum rules cannot tell us at what scale these distribution amplitudes are normalized; this scale must be low, $`\mu 1/\tau `$. The perturbative contributions vanish at $`\omega >2\epsilon _c`$. The quark-condensate contribution is the same for both $`\phi _+(\omega )`$ and $`\phi _{}(\omega )`$. Here we used the ansatz (7.17). This contribution gives a sharp peak at low $`\omega `$<sup>3</sup><sup>3</sup>3For the pion distribution amplitude, the quark-condensate contribution gives enhancements near $`x=0`$ and $`1`$.. Details of its shape are unknown, but it cannot be wider because of the restriction on the first moment (7.13). This contribution falls off quickly at larger $`\omega `$. Therefore, the distribution amplitudes at this low $`\mu `$ are only non-zero at $`\omega \overline{\mathrm{\Lambda }}`$, in accord with the expectations of the quark model.
If we neglect the quark-condensate contribution (though this is not a good idea) and also the effect of continuum subtraction<sup>4</sup><sup>4</sup>4If the same procedure is applied to the diagonal sum rule for the pion distribution amplitude, it yields the asymptotic shape., then the perturbative contributions (7.5) suggest the following simple model of distribution amplitudes:
$$\phi _+(\omega )=\frac{\omega }{\omega _0^2}e^{\omega /\omega _0},\phi _{}(\omega )=\frac{1}{\omega _0}e^{\omega /\omega _0}.$$
(7.24)
They are normalized (4.14); from the first moments (4.26) we obtain
$$\omega _0=\frac{2}{3}\overline{\mathrm{\Lambda }}.$$
(7.25)
These functions (Fig. 17) have something in common with the Wandzura–Wilczek ones (Fig. 6), but in contrast to them they fall off smoothly at large $`\omega `$.
Radiative corrections to the perturbative spectral density and the dimension-3 quark-condensate contribution for $`\phi _+`$ were calculated in . With these corrections, one can check that the correlator (7.2) satisfies the evolution equation (6.8). The resulting sum rules should be used together with the sum rules for $`F^2`$ with radiative corrections . The perturbative spectral density is
$$\begin{array}{cc}& \rho _+^{(1)}(\epsilon ,\omega )=\frac{N_c}{8\pi ^2}\omega \hfill \\ & \times \{\begin{array}{cc}x>1:\hfill & 1+C_F\frac{\alpha _s}{4\pi }[2\mathrm{log}^2\frac{\omega }{\mu }4(\mathrm{log}(x1)+1)\mathrm{log}\frac{\omega }{\mu }\hfill \\ & +2Li_2\left(\frac{1}{1x}\right)\mathrm{log}^2(x1)\hfill \\ & (2x+3)\mathrm{log}(x1)+2x\mathrm{log}x+\frac{7}{12}\pi ^2+7]\hfill \\ x<1:\hfill & 2C_F\frac{\alpha _s}{4\pi }[2(\mathrm{log}(1x)+x)\mathrm{log}\frac{\omega }{\mu }+2\mathrm{log}^2(1x)\hfill \\ & +(2x1)\mathrm{log}(1x)x]\hfill \end{array}\hfill \end{array}$$
(7.26)
where $`x=2\epsilon /\omega `$. Now the spectral density does not vanish at $`\epsilon <\omega /2`$, it is just suppressed by $`\alpha _s/(4\pi )`$. Therefore, the perturbative contribution to the sum rule (7.20) for $`\phi _+(\omega )`$ does not vanish at $`\omega >2\epsilon _c`$. It produces a radiative tail $`1/\omega `$ with the magnitude of order $`\alpha _s/(4\pi )`$.
The radiative correction to the dimension-3 quark-condensate contribution cannot be used together with a model of bilocal condensate, because radiative corrections to higher-dimensional contributions are not known. It is vital to use some model of the bilocal quark condensate in the sum rules for the distribution amplitude, because the local operator expansion produces contributions $`\delta (\omega )`$, $`\delta ^{}(\omega )`$, …, which don’t tell us much about the shape of the distribution amplitude. Therefore, the full result for this correction cannot be used in the sum rule for $`\phi _+(\omega )`$. Fortunately, the authors of demonstrated that a part of this correction is universal, and exponentialized into the Sudakov factor. We may multiply the bilocal quark-condensate contribution by this Sudakov factor.
## Acknowledgments
I am grateful to M. Neubert for collaboration on and to S.V. Mikhailov for useful discussions.
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# 1 Introduction.
## 1 Introduction.
In this note, we study an inverse scattering problem for a two-body short-range system in the presence of an external time-periodic electric field $`E(t)`$ and a time-periodic short-range potential $`V(t,x)`$ (with the same period $`T`$). For the sake of simplicity, we assume that the period $`T=1`$.
The corresponding Hamiltonian is given on $`L^2(IR^n)`$ by :
$$H(t)=\frac{1}{2}p^2E(t)x+V(t,x),$$
$`(1.1)`$
where $`p=i_x`$. When $`E(t)=0`$, the Hamiltonian $`H(t)`$ describes the dynamics of the hydrogen atom placed in a linearly polarized monochromatic electric field, or a light particle in the restricted three-body problem in which two other heavy particles are set on prescribed periodic orbits. When $`E(t)=\mathrm{cos}(2\pi t)E`$ with $`EIR^n,`$ the Hamiltonian describes the well-known AC-Stark effect in the $`E`$-direction .
In this paper, we assume that the external electric field $`E(t)`$ satisfies :
$$tE(t)L_{loc}^1(IR;IR^n),E(t+1)=E(t)a.e.$$
$`(A_1)`$
Moreover, we assume that the potential $`VC^{\mathrm{}}(IR\times IR^n)`$, is time-periodic with period $`1`$, and satisfies the following estimations :
$$\alpha IN^n,kIN,_t^k_x^\alpha V(t,x)C_{k,\alpha }<x>^{\delta \alpha },\mathrm{with}\delta >0,$$
$`(A_2)`$
where $`<x>=(1+x^2)^{\frac{1}{2}}`$. Actually, we can accommodate more singular potentials (see , , for example) and we need $`(A_2)`$ for only $`k,\alpha `$ with finite order . It is well-known that under assumptions $`(A_1)(A_2)`$, $`H(t)`$ is essentially self-adjoint on $`𝒮(IR^n)`$ the Schwartz space, . We denote $`H(t)`$ the self-adjoint realization with domain $`D(H(t))`$.
Now, let us recall some well-known results in scattering theory for time-periodic electric fields. We denote $`H_0(t)`$ the free Hamiltonian :
$$H_0(t)=\frac{1}{2}p^2E(t)x,$$
$`(1.2)`$
and let $`U_0(t,s)`$, (resp. $`U(t,s)`$) be the unitary propagator associated with $`H_0(t)`$, (resp. $`H(t)`$) (see section 2 for details).
For short-range potentials, the wave operators are defined for $`sIR`$ and $`\mathrm{\Phi }L^2(IR^n)`$ by :
$$W^\pm (s)\mathrm{\Phi }=\underset{t\pm \mathrm{}}{lim}U(s,t)U_0(t,s)\mathrm{\Phi }.$$
$`(1.3)`$
We emphasize that the short-range condition depends on the value of the mean of the external electric field :
$$E_0=_0^1E(t)𝑑t.$$
$`(1.4)`$
$``$ The case $`E_0=0`$.
By virtue of the Avron-Herbst formula (see section 2), this case falls under the category of two-body systems with time-periodic potentials and this case was studied by Kitada and Yajima (, ), Yokoyama .
We recall that for a unitary or self-adjoint operator $`U`$, $`_c(U),_{ac}(U),_{sc}(U)`$ and $`_p(U)`$ are, respectively, continuous, absolutely continuous, singular continuous and point spectral subspace of $`U`$.
We have the following result (, , ) :
###### Theorem 1
Assume that hypotheses $`(A_1),(A_2)`$ are satisfied with $`\delta >1`$ and with $`E_0=0`$.
Then : (i) the wave operators $`W^\pm (s)`$ exist for all $`sIR`$.
(ii) $`W^\pm (s+1)=W^\pm (s)`$ and $`U(s+1,s)W^\pm (s)=W^\pm (s)U_0(s+1,s)`$.
(iii) $`Ran(W^\pm (s))=_{ac}\left(U(s+1,s)\right)`$ and $`_{sc}\left(U(s+1,s)\right)=\mathrm{}`$.
(iv) the purely point spectrum $`\sigma _p(U(s+1,s))`$ is discrete outside $`\{1\}`$.
Comments.
1 - The unitary operators $`U(s+1,s)`$ are called the Floquet operators and they are mutually equivalent. The Floquet operators play a central role in the analysis of time periodic systems.
The eigenvalues of these operators are called Floquet multipliers. In , Galtbayar, Jensen and Yajima improve assertion $`(iv)`$ : for $`n=3`$ and $`\delta >2`$, $`_p\left(U(s+1,s)\right)`$ is finite dimensional.
2 - For general $`\delta >0`$, $`W^\pm (s)`$ do not exist and we have to define other wave operators. In (, ), Kitada and Yajima have constructed modified wave operators $`W_{HJ}^\pm `$ by solving an Hamilton-Jacobi equation.
$``$ The case $`E_00`$.
This case was studied by Moller : using the Avron-Herbst formula, it suffices to examine Hamiltonians with a constant external electric field, (Stark Hamiltonians) : the spectral and the scattering theory for Stark Hamiltonians are well established . In particular, a Stark Hamiltonian with a potential $`V`$ satisfying $`(A_2)`$ has no eigenvalues . The following theorem, due to Moller, is a time-periodic version of these results.
###### Theorem 2
Assume that hypotheses $`(A_1),(A_2)`$ are satisfied with $`\delta >\frac{1}{2}`$ and with $`E_00`$.
Then : (i) the Floquet operators $`U(s+1,s)`$ have purely absolutely continuous spectrum.
(ii) the wave operators $`W^\pm (s)`$ exist for all $`sIR`$ and are unitary.
(iii) $`W^\pm (s+1)=W^\pm (s)`$ and $`U(s+1,s)W^\pm (s)=W^\pm (s)U_0(s+1,s)`$.
The inverse scattering problem.
For $`sIR`$, we define the scattering operators $`S(s)=W^+(s)W^{}(s)`$. It is clear that the scattering operators $`S(s)`$ are periodic with period $`1`$.
The inverse scattering problem consists to reconstruct the perturbation $`V(s,x)`$ from the scattering operators $`S(s)`$, $`s[0,1]`$.
In this paper, we prove the following result :
###### Theorem 3
Assume that $`E(t)`$ satisfies $`(A_1)`$ and let $`V_j,j=1,2`$ be potentials satisfying $`(A_2)`$. We assume that $`\delta >1`$ (if $`E_0=0`$), $`\delta >\frac{1}{2}`$ (if $`E_00`$ and $`n3`$), $`\delta >\frac{3}{4}`$ (if $`E_00`$ and $`n=2`$). Let $`S_j(s)`$ be the corresponding scattering operators.
Then :
$$s[0,1],S_1(s)=S_2(s)V_1=V_2.$$
We prove Theorem 3 by studying the high energy limit of $`[S(s),p]`$, (Enss-Weder’s approach ). We need $`n3`$ in the case $`E_00`$ in order to use the inversion of the Radon transform on the orthogonal hyperplane to $`E_0`$. See also for a similar problem with a Stark Hamiltonian.
We can also remark that if we know the free propagator $`U_0(t,s)`$ , $`s,tIR`$, then by virtue of the following relation :
$$S(t)=U_0(t,s)S(s)U_0(s,t),$$
$`(1.5)`$
the potential $`V(t,x)`$ is uniquely reconstructed from the scattering operator $`S(s)`$ at only one initial time.
In , Yajima proves uniqueness for the case of time-periodic potential with the condition $`\delta >\frac{n}{2}+1`$ and with $`E(t)=0`$ by studying the scattering matrices in a high energy regime.
In , for a time-periodic potential that decays exponentially at infinity, Weder proves uniqueness at a fixed quasi-energy.
Note also that inverse scattering for long-range time-dependent potentials without external electric fields was studied by Weder with the Enss-Weder time-dependent method, and by Ito for time-dependent electromagnetic potentials for Dirac equations .
## 2 Proof of Theorem 3.
### 2.1 The Avron-Herbst formula.
First, let us recall some basic definitions for time-dependent Hamiltonians. Let $`\{H(t)\}_{tIR}`$ be a family of selfadjoint operators on $`L^2(IR^n)`$ such that $`𝒮(IR^n)D(H(t))`$ for all $`tIR`$.
Definition.
We call propagator a family of unitary operators on $`L^2(IR^n)`$, $`U(t,s),t,sIR`$ such that :
1 - $`U(t,s)`$ is a strongly continuous fonction of $`(t,s)IR^2`$.
2 - $`U(t,s)U(s,r)=U(t,r)`$ for all $`t,s,rIR`$.
3 - $`U(t,s)\left(𝒮(IR^n)\right)𝒮(IR^n)`$ for all $`t,sIR`$.
4 - If $`\mathrm{\Phi }𝒮(IR^n)`$, $`U(t,s)\mathrm{\Phi }`$ is continuously differentiable in $`t`$ and $`s`$ and satisfies :
$$i\frac{}{t}U(t,s)\mathrm{\Phi }=H(t)U(t,s)\mathrm{\Phi },i\frac{}{s}U(t,s)\mathrm{\Phi }=U(t,s)H(s)\mathrm{\Phi }.$$
To prove the existence and the uniqueness of the propagator for our Hamiltonians $`H(t)`$, we use a generalization of the Avron-Herbst formula close to the one given in .
In , the author gives, for $`E_00`$, a different formula which has the advantage to be time-periodic. To study our inverse scattering problem, we use here a different one, which is defined for all $`E_0`$. We emphasize that with our choice, $`c(t)`$ (see below for the definition of $`c(t)`$) is also periodic with period $`1`$; in particular $`c(t)=O(1)`$.
The basic idea is to generalize the well-known Avron-Herbst formula for a Stark Hamiltonian with a constant electric field $`E_0`$, ; if we consider the Hamiltonian $`B_0`$ on $`L^2(IR^n)`$,
$$B_0=\frac{1}{2}p^2E_0x,$$
$`(2.1)`$
we have the following formula :
$$e^{itB_0}=e^{i\frac{E_0^2}{6}t^3}e^{itE_0x}e^{i\frac{t^2}{2}E_0p}e^{it\frac{p^2}{2}}.$$
$`(2.2)`$
In the next definition, we give a similar formula for time-dependent electric fields.
Definition.
We consider the family of unitary operators $`T(t)`$, for $`tIR`$ :
$$T(t)=e^{ia(t)}e^{ib(t)x}e^{ic(t)p},$$
where :
$$b(t)=_0^t(E(s)E_0)𝑑s_0^1_0^t(E(s)E_0)𝑑s𝑑t.$$
$`(2.3)`$
$$c(t)=_0^tb(s)𝑑s.$$
$`(2.4)`$
$$a(t)=_0^t\left(\frac{1}{2}b^2(s)E_0c(s)\right)𝑑s.$$
$`(2.5)`$
###### Lemma 4
The family $`\{H_0(t)\}_{tIR}`$ has an unique propagator $`U_0(t,s)`$ defined by :
$$U_0(t,s)=T(t)e^{i(ts)B_0}T^{}(s).$$
$`(2.6)`$
Proof.
We can always assume $`s=0`$ and we make the following ansatz :
$$U_0(t,0)=e^{ia(t)}e^{ib(t)x}e^{ic(t)p}e^{itB_0}.$$
$`(2.7)`$
Since on the Schwartz space, $`U_0(t,0)`$ must satisfy :
$$i\frac{}{t}U_0(t,0)=H_0(t)U_0(t,0),$$
$`(2.8)`$
the functions $`a(t),b(t),c(t)`$ solve :
$$\dot{b}(t)=E(t)+E_0,\dot{c}(t)=b(t),\dot{a}(t)=\frac{1}{2}b^2(t)E_0c(t).$$
$`(2.9)`$
We refer to for details and for a different formula. $`\mathrm{}`$
In the same way, in order to define the propagator corresponding to the family $`\{H(t)\}`$, we consider a Stark Hamiltonian with a time-periodic potential $`V_1(t,x)`$, (we recall that $`c(t)`$ a is $`C^1`$-periodic function) :
$$B(t)=B_0+V_1(t,x)\mathrm{where}V_1(t,x)=e^{ic(t)p}V(t,x)e^{ic(t)p}=V(t,x+c(t)).$$
$`(2.10)`$
Then, $`B(t)`$ has an unique propagator $`R(t,s)`$, (see for the case $`E_0=0`$ and for the case $`E_00`$). It is easy to see that the propagator $`U(t,s)`$ for the family $`\{H(t)\}`$ is defined by :
$$U(t,s)=T(t)R(t,s)T^{}(s).$$
$`(2.11)`$
Comments.
Since the Hamiltonians $`H_0(t)`$ and $`H(t)`$ are time-periodic with period $`1`$, one has for all $`t,sIR`$ :
$$U_0(t+1,s+1)=U_0(t,s),U(t+1,s+1)=U(t,s).$$
$`(2.12)`$
Thus, the wave operators satisfy $`W^\pm (s+1)=W^\pm (s)`$.
### 2.2 The high energy limit of the scattering operators.
In this section, we study the high energy limit of the scattering operators by using the well-known Enss-Weder’s time-dependent method . This method can be used to study Hamiltonians with electric and magnetic potentials on $`L^2(IR^n)`$ , the Dirac equation , the N-body case , the Stark effect (, ), the Aharonov-Bohm effect .
In , a stationary approach, based on the same ideas, is proposed to solve scattering inverse problems for Schrödinger operators with magnetic fields or with the Aharonov-Bohm effect.
Before giving the main result of this section, we need some notation.
$``$ $`\mathrm{\Phi },\mathrm{\Psi }`$ are the Fourier transforms of functions in $`C_0^{\mathrm{}}(IR^n)`$.
$``$ $`\omega S^{n1}\mathrm{\Pi }_{E_0}`$ is fixed, where $`\mathrm{\Pi }_{E_0}`$ is the orthogonal hyperplane to $`E_0`$.
$``$ $`\mathrm{\Phi }_{\lambda ,\omega }=e^{i\sqrt{\lambda }x.\omega }\mathrm{\Phi },\mathrm{\Psi }_{\lambda ,\omega }=e^{i\sqrt{\lambda }x.\omega }\mathrm{\Psi }`$.
We have the following high energy asymptotics where $`<,>`$ is the usual scalar product in $`L^2(IR^n)`$ :
###### Proposition 5
Under the assumptions of Theorem 3, we have for all $`s[0,1]`$,
$$<[S(s),p]\mathrm{\Phi }_{\lambda ,\omega },\mathrm{\Psi }_{\lambda ,\omega }>=\lambda ^{\frac{1}{2}}<\left(_{\mathrm{}}^+\mathrm{}_xV(s,x+t\omega )dt\right)\mathrm{\Phi },\mathrm{\Psi }>+o(\lambda ^{\frac{1}{2}}).$$
Comments.
Actually, for the case $`n=2,E_00`$ and $`\delta >\frac{3}{4}`$, Proposition 5 is also valid for $`\omega S^{n1}`$ satisfying $`\omega E_0<E_0`$, (see (, ).
Then, Theorem 3 follows from Proposition 5 and the inversion of Radon transform ( and , Section 2.3).
Proof of Proposition 5.
For example, let us show Proposition 5 for the case $`E_00`$ and $`n3`$, the other cases are similar. More precisely, see for the case $`E_0=0`$, and for the case $`n=2,E_00`$, see (, Theorem 2.4) and (, Theorem 4).
Step 1.
Since $`c(t)`$ is periodic, $`c(t)=O(1)`$. Then, $`V_1(t,x)`$ is a short-range perturbation of $`B_0`$, and we can define the usual wave operators for the pair of Hamiltonians $`(B(t),B_0)`$ :
$$\mathrm{\Omega }^\pm (s)=\mathrm{s}\underset{t\pm \mathrm{}}{lim}R(s,t)e^{i(ts)B_0}.$$
$`(2.13)`$
Consider also the scattering operators $`S_1(s)=\mathrm{\Omega }^+(s)\mathrm{\Omega }^{}(s)`$. By virtue of $`(2.6)`$ and $`(2.11)`$, it is clear that :
$$S(s)=T(s)S_1(s)T^{}(s).$$
$`(2.14)`$
Using the fact that $`e^{ib(s)x}pe^{ib(s)x}=p+b(s)`$, we have :
$$[S(s),p]=[S(s),p+b(s)]=T(s)[S_1(s),p]T^{}(s).$$
$`(2.15)`$
Thus,
$$<[S(s),p]\mathrm{\Phi }_{\lambda ,\omega },\mathrm{\Psi }_{\lambda ,\omega }>=<[S_1(s),p]T^{}(s)\mathrm{\Phi }_{\lambda ,\omega },T^{}(s)\mathrm{\Psi }_{\lambda ,\omega }>.$$
$`(2.16)`$
In other hand,
$$T^{}(s)\mathrm{\Phi }_{\lambda ,\omega }=e^{i\sqrt{\lambda }x.\omega }e^{ic(s)(p+\sqrt{\lambda }\omega )}e^{ib(s)x}e^{ia(s)}\mathrm{\Phi }.$$
$`(2.17)`$
So, we obtain :
$$<[S(s),p]\mathrm{\Phi }_{\lambda ,\omega },\mathrm{\Psi }_{\lambda ,\omega }>=<[S_1(s),p]f_{\lambda ,\omega },g_{\lambda ,\omega }>,$$
$`(2.18)`$
where
$$f=e^{ic(s)p}e^{ib(s)x}\mathrm{\Phi }\mathrm{and}g=e^{ic(s)p}e^{ib(s)x}\mathrm{\Psi }.$$
$`(2.19)`$
Clearly, $`f,g`$ are the Fourier transforms of functions in $`C_0^{\mathrm{}}(IR^n)`$.
$``$ Step 2 : Modified wave operators.
Now, we follow a strategy close to for time-dependent potentials. First, let us define a free-modified dynamic $`U_D(t,s)`$ by :
$$U_D(t,s)=e^{i(ts)B_0}e^{i_0^{ts}V_1(u+s,up^{}+\frac{1}{2}u^2E_0)𝑑u},$$
$`(2.20)`$
where $`p^{}`$ is the projection of $`p`$ on the orthogonal hyperplane to $`E_0`$.
We define the modified wave operators :
$$\mathrm{\Omega }_D^\pm (s)=\mathrm{s}\underset{t\pm \mathrm{}}{lim}R(s,t)U_D(t,s).$$
$`(2.21)`$
It is clear that :
$$\mathrm{\Omega }_D^\pm (s)=\mathrm{\Omega }^\pm (s)e^{ig^\pm (s,p^{})},$$
$`(2.22)`$
where
$$g^\pm (s,p^{})=_0^\pm \mathrm{}V_1(u+s,up^{}+\frac{1}{2}u^2E_0)𝑑u.$$
$`(2.23)`$
Thus, if we set $`S_D(s)=\mathrm{\Omega }_D^+(s)\mathrm{\Omega }_D^{}(s)`$, one has :
$$S_1(s)=e^{ig^+(s,p^{})}S_D(s)e^{ig^{}(s,p^{})}$$
$`(2.24)`$
$``$ Step 3 : High energy estimates.
Denote $`\rho =min(1,\delta )`$. We have the following estimations, (the proof is exactly the same as in (, Lemma 3) for time-independent potentials).
###### Lemma 6
For $`\lambda >>1`$, we have :
$$\left(V_1(t,x)V_1(t,(ts)p^{}+\frac{1}{2}(ts)^2E_0)\right)U_D(t,s)e^{ig^\pm (s,p^{})}f_{\lambda ,\omega }$$
$`(i)`$
$$C(1+(ts)\sqrt{\lambda })^{\frac{1}{2}\rho }.$$
$$\left(R(t,s)\mathrm{\Omega }_D^\pm (s)U_D(t,s)\right)e^{ig^\pm (s,p^{})}f_{\lambda ,\omega }=O(\lambda ^{\frac{1}{2}}),\mathrm{uniformly}\mathrm{for}t,sIR.$$
$`(ii)`$
$``$ Step 4.
We denote $`F(s,\lambda ,\omega )=<[S_1(s),p]f_{\lambda ,\omega },g_{\lambda ,\omega }>`$. Using $`(2.24)`$, we have :
$$F(s,\lambda ,\omega )=<[e^{ig^+(s,p^{})}S_D(s)e^{ig^{}(s,p^{})},p]f_{\lambda ,\omega },g_{\lambda ,\omega }>$$
$$=<[S_D(s),p]e^{ig^{}(s,p^{})}f_{\lambda ,\omega },e^{ig^+(s,p^{})}g_{\lambda ,\omega }>$$
$$=<[S_D(s)1,p\sqrt{\lambda }\omega ]e^{ig^{}(s,p^{})}f_{\lambda ,\omega },e^{ig^+(s,p^{})}g_{\lambda ,\omega }>$$
$$=<(S_D(s)1)e^{ig^{}(s,p^{})}(pf)_{\lambda ,\omega },e^{ig^+(s,p^{})}g_{\lambda ,\omega }>$$
$$<(S_D(s)1)e^{ig^{}(s,p^{})}f_{\lambda ,\omega },e^{ig^+(s,p^{})}(pg)_{\lambda ,\omega }>$$
$`:=F_1(s,\lambda ,\omega )F_2(s,\lambda ,\omega )`$.
First, let us study $`F_1(s,\lambda ,\omega )`$. Writing $`S_D(s)1=(\mathrm{\Omega }_D^+(s)\mathrm{\Omega }_D^{}(s))^{}\mathrm{\Omega }_D^{}(s)`$ and using
$$\mathrm{\Omega }_D^+(s)\mathrm{\Omega }_D^{}(s)=i_{\mathrm{}}^+\mathrm{}R(s,t)\left(V_1(t,x)V_1(t,(ts)p^{}+\frac{1}{2}(ts)^2E_0)\right)U_D(t,s)𝑑t,$$
$`(2.25)`$
we obtain :
$$S_D(s)1=i_{\mathrm{}}^+\mathrm{}U_D(t,s)^{}\left(V_1(t,x)V_1(t,(ts)p^{}+\frac{1}{2}(ts)^2E_0)\right)$$
$`(2.26)`$
$$R(t,s)\mathrm{\Omega }_D^{}(s)dt.$$
Thus,
$$F_1(s,\lambda ,\omega )=i_{\mathrm{}}^+\mathrm{}<R(t,s)\mathrm{\Omega }_D^{}(s)e^{ig^{}(s,p^{})}(pf)_{\lambda ,\omega },$$
$$\left(V_1(t,x)V_1(t,(ts)p^{}+\frac{1}{2}(ts)^2E_0)\right)U_D(t,s)e^{ig^+(s,p^{})}g_{\lambda ,\omega }>dt$$
$$=i_{\mathrm{}}^+\mathrm{}<U_D(t,s)e^{ig^{}(s,p^{})}(pf)_{\lambda ,\omega },$$
$$\left(V_1(t,x)V_1(t,(ts)p^{}+\frac{1}{2}(ts)^2E_0)\right)U_D(t,s)e^{ig^+(s,p^{})}g_{\lambda ,\omega }>dt$$
$`+R_1(s,\lambda ,\omega ),`$
where :
$$R_1(s,\lambda ,\omega )=i_{\mathrm{}}^+\mathrm{}<\left(R(t,s)\mathrm{\Omega }_D^{}(s)U_D(t,s)\right)e^{ig^{}(s,p^{})}(pf)_{\lambda ,\omega },$$
$`(2.27)`$
$$\left(V_1(t,x)V_1(t,(ts)p^{}+\frac{1}{2}(ts)^2E_0)\right)U_D(t,s)e^{ig^+(s,p^{})}g_{\lambda ,\omega }>dt.$$
By Lemma 6, it is clear that $`R_1(s,\lambda ,\omega )=O(\lambda ^1)`$. Thus, writing $`t=\frac{\tau }{\sqrt{\lambda }}+s`$, we obtain :
$$F_1(s,\lambda ,\omega )=\frac{i}{\sqrt{\lambda }}_{\mathrm{}}^+\mathrm{}<U_D(\frac{\tau }{\sqrt{\lambda }}+s,s)e^{ig^{}(s,p^{})}(pf)_{\lambda ,\omega },$$
$`(2.28)`$
$$\left(V_1(\frac{\tau }{\sqrt{\lambda }}+s,x)V_1(\frac{\tau }{\sqrt{\lambda }}+s,\frac{\tau }{\sqrt{\lambda }}p^{}+\frac{\tau ^2}{2\lambda }E_0)\right)$$
$$U_D(\frac{\tau }{\sqrt{\lambda }}+s,s)e^{ig^+(s,p^{})}g_{\lambda ,\omega }>d\tau +O(\lambda ^1).$$
Denote by $`f_1(\tau ,s,\lambda ,\omega )`$ the integrand of the (R.H.S) of $`(2.28)`$. By Lemma 6 (i),
$$f_1(\tau ,s,\lambda ,\omega )C(1+\tau )^{\frac{1}{2}\rho }.$$
$`(2.29)`$
So, by Lebesgue’s theorem, to obtain the asymptotics of $`F_1(s,\lambda ,\omega )`$, it suffices to determine $`\underset{\lambda +\mathrm{}}{lim}f_1(\tau ,s,\lambda ,\omega )`$.
Let us denote :
$$U^\pm (t,s,p^{})=e^{i_t^\pm \mathrm{}V_1(u+s,up^{}+\frac{1}{2}u^2E_0)𝑑u}.$$
$`(2.30)`$
We have :
$$f_1(\tau ,s,\lambda ,\omega )=<e^{i\frac{\tau }{\sqrt{\lambda }}B_0}U^{}(\frac{\tau }{\sqrt{\lambda }},s,p^{})(pf)_{\lambda ,\omega },$$
$`(2.31)`$
$$\left(V_1(\frac{\tau }{\sqrt{\lambda }}+s,x)V_1(\frac{\tau }{\sqrt{\lambda }}+s,\frac{\tau }{\sqrt{\lambda }}p^{}+\frac{\tau ^2}{2\lambda }E_0)\right)e^{i\frac{\tau }{\sqrt{\lambda }}B_0}U^+(\frac{\tau }{\sqrt{\lambda }},s,p^{})g_{\lambda ,\omega }>.$$
Using the Avron-Herbst formula $`(2.2)`$, we deduce that :
$$f_1(\tau ,s,\lambda ,\omega )=<e^{i\frac{\tau }{2\sqrt{\lambda }}p^2}U^{}(\frac{\tau }{\sqrt{\lambda }},s,p^{})(pf)_{\lambda ,\omega },$$
$`(2.32)`$
$$\left(V_1(\frac{\tau }{\sqrt{\lambda }}+s,x+\frac{\tau ^2}{2\lambda }E_0)V_1(\frac{\tau }{\sqrt{\lambda }}+s,\frac{\tau }{\sqrt{\lambda }}p^{}+\frac{\tau ^2}{2\lambda }E_0)\right)e^{i\frac{\tau }{2\sqrt{\lambda }}p^2}U^+(\frac{\tau }{\sqrt{\lambda }},s,p^{})g_{\lambda ,\omega }>.$$
Then, we obtain :
$$f_1(\tau ,s,\lambda ,\omega )=<e^{i\frac{\tau }{2\sqrt{\lambda }}(p+\sqrt{\lambda }\omega )^2}U^{}(\frac{\tau }{\sqrt{\lambda }},s,p^{}+\sqrt{\lambda }\omega )pf,$$
$`(2.33)`$
$$\left(V_1(\frac{\tau }{\sqrt{\lambda }}+s,x+\frac{\tau ^2}{2\lambda }E_0)V_1(\frac{\tau }{\sqrt{\lambda }}+s,\frac{\tau }{\sqrt{\lambda }}(p^{}+\sqrt{\lambda }\omega )+\frac{\tau ^2}{2\lambda }E_0)\right)$$
$$e^{i\frac{\tau }{2\sqrt{\lambda }}(p+\sqrt{\lambda }\omega )^2}U^+(\frac{\tau }{\sqrt{\lambda }},s,p^{}+\sqrt{\lambda }\omega )g>.$$
Since
$$e^{i\frac{\tau }{2\sqrt{\lambda }}(p+\sqrt{\lambda }\omega )^2}=e^{i\frac{\tau \sqrt{\lambda }}{2}}e^{i\tau \omega .p}e^{i\frac{\tau }{2\sqrt{\lambda }}p^2},$$
$`(2.34)`$
we have
$$f_1(\tau ,s,\lambda ,\omega )=<e^{i\frac{\tau }{2\sqrt{\lambda }}p^2}U^{}(\frac{\tau }{\sqrt{\lambda }},s,p^{}+\sqrt{\lambda }\omega )pf,$$
$`(2.35)`$
$$\left(V_1(\frac{\tau }{\sqrt{\lambda }}+s,x+\tau \omega +\frac{\tau ^2}{2\lambda }E_0)V_1(\frac{\tau }{\sqrt{\lambda }}+s,\frac{\tau }{\sqrt{\lambda }}(p^{}+\sqrt{\lambda }\omega )+\frac{\tau ^2}{2\lambda }E_0)\right)$$
$$e^{i\frac{\tau }{2\sqrt{\lambda }}p^2}U^+(\frac{\tau }{\sqrt{\lambda }},s,p^{}+\sqrt{\lambda }\omega )g>.$$
Since $`V_1(u+s,u(p^{}+\sqrt{\lambda }\omega )+\frac{1}{2}u^2E_0))C(u^2+1)^\delta L^1(IR^+,du)`$, it is easy to show (using Lebesgue’s theorem again) that :
$$s\underset{\lambda +\mathrm{}}{lim}U^\pm (\frac{\tau }{\sqrt{\lambda }},s,p^{}+\sqrt{\lambda }\omega )=1.$$
$`(2.38)`$
Then,
$$\underset{\lambda +\mathrm{}}{lim}f_1(\tau ,s,\lambda ,\omega )=<pf,(V_1(s,x+\tau \omega )V_1(s,\tau \omega ))g>.$$
$`(2.39)`$
So, we have obtained :
$$F_1(s,\lambda ,\omega )=\frac{i}{\sqrt{\lambda }}<pf,\left(_{\mathrm{}}^+\mathrm{}(V_1(s,x+\tau \omega )V_1(s,\tau \omega ))𝑑\tau \right)g>+o(\frac{1}{\sqrt{\lambda }}).$$
$`(2.40)`$
In the same way, we obtain
$$F_2(s,\lambda ,\omega )=\frac{i}{\sqrt{\lambda }}<f,\left(_{\mathrm{}}^+\mathrm{}(V_1(s,x+\tau \omega )V_1(s,\tau \omega ))𝑑\tau \right)pg>+o(\frac{1}{\sqrt{\lambda }}),$$
$`(2.41)`$
so
$$F(s,\lambda ,\omega )=F_1(s,\lambda ,\omega )F_2(s,\lambda ,\omega )$$
$`(2.42)`$
$$=\frac{1}{\sqrt{\lambda }}<f,\left(_{\mathrm{}}^+\mathrm{}_xV_1(s,x+\tau \omega )d\tau \right)g>+o(\frac{1}{\sqrt{\lambda }}).$$
$`(2.43)`$
Using $`(2.19)`$ and $`_xV(s,x+\tau \omega )=e^{ic(s)p}_xV_1(s,x+\tau \omega )e^{ic(s)p}`$, we obtain :
$$F(s,\lambda ,\omega )=\frac{1}{\sqrt{\lambda }}<\mathrm{\Phi },\left(_{\mathrm{}}^+\mathrm{}_xV(s,x+\tau \omega )d\tau \right)\mathrm{\Psi }>+o(\frac{1}{\sqrt{\lambda }}).\mathrm{}$$
$`(2.44)`$
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# Stability analysis of dynamical regimes in nonlinear systems with discrete symmetries
## I Introduction
Different dynamical regimes in a mechanical system with discrete-symmetry group $`G_0`$ can be classified by subgroups $`G_jG_0`$ of this group DAN1 ; DAN2 ; PhysD . Actually, we can find an invariant manifold corresponding to each subgroup $`G_j`$ and decompose it into the basis vectors of the irreducible representations of the group $`G_0`$. As a result of this procedure, we obtain a *bush of modes* (see above cited papers) which can be considered as a certain physical object in geometrical, as well as dynamical sense. The mode structure of a given bush is fully determined by its symmetry group $`G_j`$ and is independent of the specific type of interparticle interactions in the system. In Hamiltonian systems, bushes of modes represent dynamical objects in which the energy of initial excitation turns out to be “trapped” (this is a phenomenon of energy localization in the modal space). The number of modes belonging to a given bush (the bush dimension) does not change in time, while amplitudes of the modes do change, and we can find dynamical equations determining their evolution.
Being an *exact* nonlinear excitation, in the considered mechanical system, each bush possesses its own domain of *stability* depending on the value of its mode amplitudes. Beyond the stability threshold a phenomenon similar to the parametric resonance occurs, the bush loses its stability and transforms into another bush of higher dimension. This process is accompanied by spontaneous lowering of the bush symmetry: $`G_j\stackrel{~}{G_j}`$, where $`\stackrel{~}{G_j}G_j`$.
The concept of bushes of modes was introduced in DAN1 ; DAN2 , the detailed theory of these dynamical objects was developed in PhysD . Low-dimensional bushes in mechanical systems with various kinds of symmetry and structures were studied in DAN1 ; DAN2 ; PhysD ; IntJ ; ENOC ; Octa ; C60 ; FPU1 ; FPU2 . The problem of bush stability was discussed in PhysD ; Octa ; FPU1 ; FPU2 . Two last papers are devoted to the vibrational bushes in the Fermi-Pasta-Ulam (FPU) chains.
Note that dynamical objects equivalent to the bushes of modes were recently discussed for the monoatomic chains in the papers of different authors PR ; BR ; Shin ; AntiFPU . Let us emphasize that group-theoretical methods developed in our papers DAN1 ; DAN2 ; PhysD can be applied efficiently not only to the monoatomic chains (as was illustrated in FPU1 ; FPU2 ), but to *all* other physical systems with discrete symmetry groups (see, DAN1 ; DAN2 ; PhysD ; IntJ ; ENOC ; Octa ; C60 ).
In this paper, we present a theorem which can simplify essentially the stability analysis of the bushes of modes in complex systems with many degrees of freedom. The usefulness of this theorem is illustrated with the example of nonlinear chains with a large number of particles. Note that the simplification of the stability analysis in such systems actually originate from the well-known Wigner theorem about the block-diagonalization of the matrix commuting with all matrices of a representation of a given symmetry group.
In Sec. II, we start with the simplest examples for introducing basic concepts and ideas. In Sec. III, we present a general theorem about invariance of the dynamical equations linearized near a given bush with respect to the bush symmetry group. In Sec. IV, we prove a theorem which turns out to be very useful for splitting the above mentioned linearized dynamical equations for $`N`$-particle monoatomic chains. Some results on the bush stability in the FPU chains are discussed in Sec. V.
## II Some simple examples
### II.1 FPU-chains and their symmetry
We consider *longitudinal* vibrations of $`N`$-particle chains of identical masses ($`m=1`$) and identical springs connecting neighboring particles. Let $`x_i(t)`$ be the displacement of the $`i`$-th particle ($`i=1,2,\mathrm{},N`$) from its equilibrium position at a given instance $`t`$. Dynamical equations of such mechanical system (FPU-chain) can be written as follows:
$$\ddot{x}_i=f(x_{i+1}x_i)f(x_ix_{i1}).$$
(1)
The nonlinear force $`f(x)`$ depends on the deformation $`x`$ of the spring as $`f(x)=x+x^2`$ and $`f(x)=x+x^3`$ for the FPU-$`\alpha `$ and FPU-$`\beta `$ chains, respectively. We assume the periodic boundary conditions
$$x_0(t)x_N(t),x_{N+1}(t)x_1(t)$$
(2)
to be valid. Let us also introduce the “configuration vector” $`𝑿(t)`$ which is the $`N`$-dimensional vector describing all the displacements of the individual particles at the moment $`t`$:
$$𝑿(t)=\{x_1(t),x_2(t),\mathrm{},x_N(t)\}.$$
(3)
In the equilibrium state, a given chain is invariant under the action of the operator $`\widehat{a}`$ which shifts the chain by the lattice spacing $`a`$. This operator generates the translational group
$$T_N=\{\widehat{e},\widehat{a},\widehat{a}^2,\mathrm{},\widehat{a}^{N1}\},\widehat{a}^N=\widehat{e},$$
(4)
where $`\widehat{e}`$ is the identity element and $`N`$ is the order of the cyclic group $`T_N`$. The operator $`\widehat{a}`$ induces the cyclic permutation of all particles of the chain and, therefore, it acts on the “configuration vector” $`𝑿(t)`$ as follows:
$$\widehat{a}𝑿(t)\widehat{a}\{x_1(t),x_2(t),\mathrm{},x_{N1}(t),x_N(t)\}=\{x_N(t),x_1(t),x_2(t),\mathrm{},x_{N1}(t)\}.$$
The full symmetry group of the monoatomic chain contains also the inversion $`\widehat{ı}`$, with respect to the center of the chain, which acts on the vector $`𝑿(t)`$ in the following manner:
$$\widehat{ı}𝑿(t)\widehat{ı}\{x_1(t),x_2(t),\mathrm{},x_{N1}(t),x_N(t)\}=\{x_N(t),x_{N1}(t),\mathrm{},x_2(t),x_1(t)\}.$$
The complete set of all products $`\widehat{a}^k\widehat{ı}`$ of the pure translations $`\widehat{a}^k`$ ($`k=0,1,2,\mathrm{},N1`$) with the inversion $`\widehat{ı}`$ forms the so-called dihedral group $`D_N`$ which can be written as the direct sum of two cosets $`T_N`$ and $`T_N\widehat{ı}`$:
$$D_N=T_NT_N\widehat{ı}.$$
(5)
The dihedral group is a non-Abelian group induced by two generators ($`\widehat{a}`$ and $`\widehat{ı}`$) with the following generating relations
$$\widehat{a}^N=\widehat{e},\widehat{ı}^2=\widehat{e},\widehat{ı}\widehat{a}=\widehat{a}^1\widehat{ı}.$$
(6)
We will consider different vibrational regimes in the FPU chains, which can be determined by the *specific forms* of the configuration vector. Each of these regimes depends on $`m`$ independent parameters ($`mN`$) and this number is the *dimension* of the given regime.
The simplest case of one-dimensional vibrational regimes represents the so-called $`\pi `$-mode (zone boundary mode) <sup>1</sup><sup>1</sup>1The dots in $`𝑿(t)`$ denote that the displacement fragment which is given explicitly must be repeated several times to form the full displacement pattern corresponding to the given bush.:
$$𝑿(t)=\{A(t),A(t)|A(t),A(t)|A(t),A(t)|\mathrm{}\},$$
(7)
where $`A(t)`$ is a certain function of $`t`$. In our terminology, this is the one-dimensional bush B$`[\widehat{a}^2,\widehat{ı}]`$ (see below about notation of bushes of modes).
The vector
$$𝑿(t)=\{0,A(t),B(t),0,B(t),A(t)|\mathrm{}\},$$
(8)
represents a two-dimensional vibrational regime that is determined by two time-dependent functions $`A(t)`$ and $`B(t)`$. This is the two-dimensional bush B$`[\widehat{a}^6,\widehat{a}\widehat{ı}]`$ (see FPU2 ).
In general, for $`m`$-dimensional vibrational regime, we write $`𝑿(t)=𝑪(t)`$, where the $`N`$-dimensional vector $`𝑪(t)`$ depends on $`m`$ time-dependent functions only. Each specific dynamical regime $`𝑪(t)`$, being an invariant manifold, possesses its own symmetry group that is a subgroup of the parent symmetry group $`G_0=D_N`$ of the chain in equilibrium.
### II.2 FPU-$`\alpha `$ chain with $`N=4`$ particles: existence of the bush B$`[\widehat{a}^2,\widehat{ı}]`$
Let us consider the above discussed equations for the simplest case $`N=4`$. Dynamical equations (1) read:
$$\begin{array}{c}\ddot{x}_1=f(x_2x_1)f(x_1x_4),\hfill \\ \ddot{x}_2=f(x_3x_2)f(x_2x_1),\hfill \\ \ddot{x}_3=f(x_4x_3)f(x_3x_2),\hfill \\ \ddot{x}_4=f(x_1x_4)f(x_4x_3).\hfill \end{array}$$
(9)
The symmetry group $`G_0`$ in the equilibrium state reads:
$$G_0=D_4=[\widehat{a},\widehat{ı}]=\{\widehat{e},\widehat{a},\widehat{a}^2,\widehat{a}^3,\widehat{ı},\widehat{a}\widehat{ı},\widehat{a}^2\widehat{ı},\widehat{a}^3\widehat{ı}\}.$$
Hereafter, we write generators of any symmetry group in square brackets, while all its elements (if it is necessary) are given in curly brackets.
The operators $`\widehat{a}`$ and $`\widehat{ı}`$ act on the configuration vector $`𝑿=\{x_1,x_2,x_3,x_4\}`$ as follows
$$\widehat{a}𝑿=\{x_4,x_1,x_2,x_3\},\widehat{ı}𝑿=\{x_4,x_3,x_2,x_1\}.$$
Therefore, we can associate the following matrices $`\mathrm{M}(\widehat{a})`$ and $`\mathrm{M}(\widehat{ı})`$ of the *mechanical representation* with these generators:
$$\widehat{a}\mathrm{M}(\widehat{a})=\left(\begin{array}{cccc}\hfill 0& \hfill 0& \hfill 0& \hfill 1\\ \hfill 1& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 1& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 1& \hfill 0\end{array}\right),\widehat{ı}\mathrm{M}(\widehat{ı})=\left(\begin{array}{cccc}\hfill 0& \hfill 0& \hfill 0& \hfill 1\\ \hfill 0& \hfill 0& \hfill 1& \hfill 0\\ \hfill 0& \hfill 1& \hfill 0& \hfill 0\\ \hfill 1& \hfill 0& \hfill 0& \hfill 0\end{array}\right).$$
(10)
Their action on the configuration vector $`𝑿`$ is equivalent to that of the operators $`\widehat{a}`$ and $`\widehat{ı}`$, respectively.
Let us now make the transformations of variables in the system (9) according to the action of the matrices (10), i.e.
$$\mathrm{M}(\widehat{a}):x_1x_4,x_2x_1,x_3x_2,x_4x_3,$$
(11)
$$\mathrm{M}(\widehat{ı}):x_1x_4,x_2x_3,x_3x_2,x_4x_1.$$
(12)
It is easy to check that both transformations (11) and (12) produce systems of equations which are *equivalent* to the system (9). Moreover, these transformations act on the individual equations $`u_j`$ ($`j=1,2,3,4`$) of the system (9) exactly as on the components $`x_j`$ ($`j=1,2,3,4`$) of the configuration vector $`𝑿`$. For example, for the operator $`\widehat{ı}`$ (or matrix $`\mathrm{M}(\widehat{ı})`$) we have
$$\widehat{ı}u_1=u_4,\widehat{ı}u_2=u_3,\widehat{ı}u_3=u_2,\widehat{ı}u_4=u_1.$$
It is obvious that a certain transposition of these equations multiplied by $`\pm 1`$ will indeed, produce a system fully identical to the original system (9). Thus, we are convinced that the symmetry group $`G_0=D_4`$ of our chain in equilibrium turns out to be the symmetry group (the group of invariance) of the dynamical equations of this mechanical system.
Let us now consider the vibrational regime (7), i.e. $`\pi `$-mode, and check that it represents an invariant manifold for the dynamical system (9). Substituting $`x_1(t)=x_3(t)=A(t)`$, $`x_2(t)=x_4(t)=A(t)`$ into (9), we reduce these equations to one and the same equation of the form
$$\ddot{A}=f(2A)f(2A).$$
(13)
In the case of the FPU-$`\alpha `$ model this equation turns out to be the equation of harmonic oscillator (for the FPU-$`\beta `$ model it reduces to the Duffing equation). Indeed, for the FPU-$`\alpha `$ chain, we obtain from the Eq.(13):
$$\ddot{A}+4A=0.$$
(14)
Using, for simplicity, the initial condition $`A(0)=C_0`$, $`\dot{A}(0)=0`$, we get the following solution to Eq. (14):
$$A(t)=C_0\mathrm{cos}(2t).$$
(15)
Thus, the one-dimensional bush B$`[\widehat{a}^2,\widehat{ı}]`$ (or $`\pi `$-mode) (7) for the FPU-$`\alpha `$ chain, represents purely harmonic dynamical regime
$$𝑿(t)=C_0\{\mathrm{cos}(2t),\mathrm{cos}(2t)|\mathrm{cos}(2t),\mathrm{cos}(2t)\}.$$
(16)
On the other hand, the invariant manifold $`𝑿(t)=\{A(t),A(t),A(t),A(t)\}`$, corresponding to the bush B$`[\widehat{a}^2,\widehat{ı}]`$, can be obtained with the aid of the group-theoretical methods only, *without* consideration of the dynamical equations (9). Let us discuss this point in more detail.
At an arbitrary instant $`t`$, the displacement pattern $`\{A(t),A(t),A(t),A(t)\}`$ possesses its own symmetry group $`G=D_2G_0=D_4`$. Indeed, this pattern is conserved under inversion ($`\widehat{ı}`$) and under shifting all particles by $`2a`$. The latter procedure can be considered as a result of the action on the chain by the operator $`\widehat{a}^2`$. These two symmetry elements ($`\widehat{a}^2`$ and $`\widehat{ı}`$) determine the dihedral group $`D_2`$ which is a subgroup of order two of the original group $`G_0=D_4`$ <sup>2</sup><sup>2</sup>2Let us recall that the order $`m`$ of the subgroup $`G`$ in the group $`G_0`$ is determined by the equation $`m=G_0/G`$, where $`G_0`$ and $`G`$ are numbers of elements in the groups $`G_0`$ and $`G`$, respectively..
It is obvious, that the old element $`\widehat{a}`$ of the group $`G_0`$, describing the chain in equilibrium, *does not survive* in the vibrational state described by the pattern (7). Note that this element ($`\widehat{a}`$) transforms the regime (7) into its equivalent (but different!) form $`\{A(t),A(t),A(t),A(t)\}`$. In the present paper, we will not discuss different equivalent forms of bushes of modes (a detailed consideration of this problem can be found in FPU2 ). Thus, we encounter the *reduction* of symmetry $`G_0=D_4G=D_2`$ when we pass from the equilibrium state to the vibrational state (7) for the considered mechanical system.
The dynamical regime (7) represents the one-dimensional bush consisting of only one mode ($`\pi `$-mode). We will denote it as B$`[G]`$=B$`[\widehat{a}^2,\widehat{ı}]:\{A,A,A,A\}`$. In square brackets, the group of the bush symmetry is indicated by listing its generators ($`\widehat{a}^2`$ and $`\widehat{ı}`$, in our case), while the characteristic fragment of the bush displacement pattern is presented next to the colon. The bush symmetry group $`G`$ fully determines the form (displacement pattern) of the bush B$`[G]`$ (see, for example, PhysD ; FPU2 ). Indeed, in the case of the bush B$`[\widehat{a}^2,\widehat{ı}]`$, it is easy to show that this form, $`𝑿=\{A(t),A(t),A(t),A(t)\}`$, can be obtained as the general solution to the following linear algebraic equation representing the invariance of the configuration vector $`𝑿`$: $`\widehat{g}_1𝑿=𝑿`$, $`\widehat{g}_2𝑿=𝑿`$, where $`\widehat{g}_1=\widehat{a}^2`$ and $`\widehat{g}_2=\widehat{ı}`$ are the generators of the group $`G`$. In our previous papers we often write these invariance conditions for the bush B$`[G]`$ in the form
$$\widehat{G}𝑿=𝑿.$$
(17)
It is very essential, that the invariant vector $`𝑿(t)`$, which was found in such *geometrical* (group-theoretical) manner, turns out to be an invariant manifold for the considered dynamical system PhysD . Thus, we can obtain the symmetry-determined invariant manifolds (bushes of modes) without any information on interparticle interactions in the mechanical system.
### II.3 FPU-$`\alpha `$ chain with $`N=4`$ particles: stability of the bush B$`[\widehat{a}^2,\widehat{ı}]`$
We now turn to the question of the stability of the bush B$`[\widehat{a}^2,\widehat{ı}]`$, representing a *periodic* vibrational regime $`𝑿=\{A(t),A(t),A(t),A(t)\}`$, with $`A(t)=C_0\mathrm{cos}(2t)`$. According to the conventional prescription, we must linearize the dynamical system (9) in the *infinitesimal vicinity* of the given bush and then study the obtained system. For this goal, let us write
$$𝑿(t)=𝑪(t)+𝜹(t),$$
(18)
where $`𝑪=\{A(t),A(t),A(t),A(t)\}`$ represents our bush, while $`𝜹(t)=\{\delta _1(t),\delta _2(t),\delta _3(t),\delta _4(t)\}`$ is an infinitesimal vector. Substituting (18) into Eqs. (9) and neglecting all terms nonlinear in $`\delta _j(t)`$, we obtain the following linearized equations for the FPU-$`\alpha `$ model:
$$\begin{array}{c}\ddot{\delta }_1=[\delta _22\delta _1+\delta _4]4A(t)[\delta _2\delta _4],\hfill \\ \ddot{\delta }_2=[\delta _32\delta _2+\delta _1]+4A(t)[\delta _3\delta _1],\hfill \\ \ddot{\delta }_3=[\delta _42\delta _3+\delta _2]4A(t)[\delta _4\delta _2],\hfill \\ \ddot{\delta }_4=[\delta _12\delta _4+\delta _3]+4A(t)[\delta _1\delta _3].\hfill \end{array}$$
(19)
The last system of equations can be written in the form
$$\ddot{𝜹}=\mathrm{J}(t)𝜹,$$
(20)
where $`\mathrm{J}(t)`$ is the *Jacobi* matrix for the system (9) calculated by the substitution of the vector $`𝑿=\{A(t),A(t),A(t),A(t)\}`$. This matrix can be presented as follows:
$$\mathrm{J}(t)=\mathrm{L}+4A(t)\mathrm{M},$$
(21)
where
$$\mathrm{L}=\left(\begin{array}{cccc}\hfill 2& \hfill 1& \hfill 0& \hfill 1\\ \hfill 1& \hfill 2& \hfill 1& \hfill 0\\ \hfill 0& \hfill 1& \hfill 2& \hfill 1\\ \hfill 1& \hfill 0& \hfill 1& \hfill 2\end{array}\right),\mathrm{M}=\left(\begin{array}{cccccc}\hfill 0& \hfill 1& \hfill 0& \hfill 1& & \\ \hfill 1& \hfill 0& \hfill 1& \hfill 0& & \\ \hfill 0& \hfill 1& \hfill 0& \hfill 1& & \\ \hfill 1& \hfill 0& \hfill 1& \hfill 0& & \end{array}\right)$$
(22)
are two time-independent symmetric matrices.
It easy to check that matrices $`\mathrm{L}`$ and $`\mathrm{M}`$ commute with each other: $`\mathrm{L}\mathrm{M}=\mathrm{M}\mathrm{L}`$. Therefore, there exists a time-independent orthogonal matrix $`\mathrm{S}`$ that transforms the both matrices $`\mathrm{L}`$ and $`\mathrm{M}`$ to the diagonal form: $`\stackrel{~}{\mathrm{S}}\mathrm{L}\mathrm{S}=\mathrm{L}_{dia}`$, $`\stackrel{~}{\mathrm{S}}\mathrm{M}\mathrm{S}=\mathrm{M}_{dia}`$ (here $`\stackrel{~}{\mathrm{S}}`$ is the transposed matrix with respect to $`\mathrm{S}`$). In turn, it means that the Jacobi matrix $`\mathrm{J}(t)`$ can be diagonalized at *any time* $`t`$ by one and the same time-independent matrix $`\mathrm{S}`$. Therefore, our linearized system (20) for the considered bush B$`[\widehat{a}^2,\widehat{ı}]`$ can be decomposed into four *independent* differential equations.
Let us discuss how the above matrix $`\mathrm{S}`$ can be obtain with the aid of the theory of irreducible representations of the symmetry group $`G`$ (in our case $`G=D_2`$).
In Sec. IV, we will consider a general method for obtaining the matrix $`\mathrm{S}`$ which reduces the Jacobi matrix $`\mathrm{J}(t)`$ to a block-diagonal form. This method uses the basis vectors of irreducible representations of the group $`G`$, constructed in the mechanical space of the considered dynamical system. In our simplest case of the monoatomic chain with $`N=4`$ particles, this method leads to the following result
$$\mathrm{S}=\frac{1}{2}\left(\begin{array}{cccc}\hfill 1& \hfill 1& \hfill 1& \hfill 1\\ \hfill 1& \hfill 1& \hfill 1& \hfill 1\\ \hfill 1& \hfill 1& \hfill 1& \hfill 1\\ \hfill 1& \hfill 1& \hfill 1& \hfill 1\end{array}\right).$$
(23)
The rows of the matrix $`\mathrm{S}`$ from (23) are simply the characters of four one-dimensional irreducible representations (irreps) – $`\mathrm{\Gamma }_1`$, $`\mathrm{\Gamma }_2`$, $`\mathrm{\Gamma }_3`$, $`\mathrm{\Gamma }_4`$ – of the Abelian group $`D_2`$, because each of these irreps is contained once in the decomposition of the mechanical representation of the group $`G=D_2`$. Introducing new variables $`𝒚=\{y_1,y_2,y_3,y_4\}`$ instead of the old variables $`𝜹=\{\delta _1(t),\delta _2(t),\delta _3(t),\delta _4(t)\}`$ by the equation $`𝒚=\mathrm{S}𝜹`$ with $`\mathrm{S}`$ from (23), we arrive at the full splitting of the linearized equations (19) for the FPU-$`\alpha `$ model:
$`\ddot{y}_1=0,`$ (24a)
$`\ddot{y}_2=2[1+4A(t)]y_2,`$ (24b)
$`\ddot{y}_3=4y_3,`$ (24c)
$`\ddot{y}_4=2[14A(t)]y_4,`$ (24d)
where $`A(t)=C_0\mathrm{cos}(2t)`$.
With the aid of Eqs. (24), we can find the stability threshold in $`C_0`$ for loss of stability of the one-dimensional bush B$`[\widehat{a}^2,\widehat{ı}]`$. Indeed, according to Eqs. (24), the variables $`y_j(t)`$ ($`j=1,2,3,4`$) are independent from each other, and we can consider them in turn. Eq. (24a) for $`y_1(t)`$ describes the uniform motion of the center of masses of our chain, since it follows from the equations $`𝒚=\mathrm{S}𝜹`$ that $`y_1(t)(\delta _1(t)+\delta _2(t)+\delta _3(t)+\delta _4(t))`$. Therefore, considering *vibrational* regimes only, we may assume $`y_1(t)0`$.
If only $`y_3(t)`$ appears in the solution to the system (24), i.e. if $`y_1(t)=0`$, $`y_2(t)=0`$, $`y_4(t)=0`$, then we have from the equation $`𝜹=\stackrel{~}{\mathrm{S}}𝒚`$ (note that $`\mathrm{S}`$ is the orthogonal matrix and, therefore, $`\mathrm{S}^1=\stackrel{~}{\mathrm{S}}`$): $`𝜹(t)=\{y_3(t),y_3(t),y_3(t),y_3(t)\}`$, where $`y_3(t)\mathrm{cos}(2t)`$. This solution leads only to deviations “along” the bush $`𝑿(t)=C_0\{\mathrm{cos}(2t),\mathrm{cos}(2t),\mathrm{cos}(2t),\mathrm{cos}(2t)\}`$ and does not signify instability.
Since $`A(t)=C_0\mathrm{cos}(2t)`$, Eq. (24b) reads $`\ddot{y}_2+[2+8C_0\mathrm{cos}(2t)]y_2=0`$ and can be transformed to the standard form of the Mathieu equation, as well as Eq. (24d). Therefore, the stability threshold of the considered bush B$`[\widehat{a}^2,\widehat{ı}]`$ for $`N=4`$ can be determined directly from the well-known diagram of the regions of stable and unstable motion of the Mathieu equation. In such a way we can find that critical value $`C_c`$ for the amplitude $`C_0`$ of the given bush for which it loses its stability is $`C_c`$ =$`0.303`$.
In conclusion, let us focus on the point that turns out to be very important for proving the general theorem in Sec. III. The system (19) was obtained by linearizing the original system (9), near the dynamical regime $`𝑪(t)=\{A(t),A(t),A(t),A(t)\}`$, and Eqs. (9) are invariant with respect to the parent group $`G_0=[\widehat{a},\widehat{ı}]`$. Despite this fact, Eqs. (19) are invariant only with respect to its *subgroup* $`G=\{\widehat{e},\widehat{a}^2,\widehat{ı},\widehat{a}^2\widehat{ı}\}G_0=\{\widehat{e},\widehat{a},\widehat{a}^2,\widehat{a}^3,\widehat{ı},\widehat{a}\widehat{ı},\widehat{a}^2\widehat{ı},\widehat{a}^3\widehat{ı}\}`$: the element $`\widehat{a}G_0`$ (as well as $`\widehat{a}^3`$, $`\widehat{a}\widehat{ı}`$, $`\widehat{a}^3\widehat{ı}`$) does not survive as a result of the symmetry reduction $`G_0G`$. Indeed, acting on Eqs. (19) by the operator $`\widehat{g}=\widehat{a}`$, which transposes variables $`\delta _j`$ as follows
$$\delta _1\delta _4,\delta _2\delta _1,\delta _3\delta _2,\delta _4\delta _3,$$
(25)
we obtain the equations:
$$\begin{array}{c}\ddot{\delta }_4=[\delta _12\delta _4+\delta _3]4A(t)[\delta _1\delta _3],\hfill \\ \ddot{\delta }_1=[\delta _22\delta _1+\delta _4]+4A(t)[\delta _2\delta _4],\hfill \\ \ddot{\delta }_2=[\delta _32\delta _2+\delta _1]4A(t)[\delta _3\delta _1],\hfill \\ \ddot{\delta }_3=[\delta _42\delta _3+\delta _2]+4A(t)[\delta _4\delta _2].\hfill \end{array}$$
(26)
Obviously, this system is not equivalent to the system (19)! (The equivalence between (19) and (26) can be restored, if, besides cyclic permutation (25) in Eqs. (19), we add the artificial transformation $`A(t)A(t)`$).
What is the source of this phenomenon? The original nonlinear dynamical system, which can be written as $`\ddot{𝑿}=𝑭(𝑿)`$, is invariant under the action of the operator $`\widehat{g}=\widehat{a}`$. Being linearized, by the substitution $`𝑿(t)=𝑪(t)+𝜹(t)`$ and neglecting all the nonlinear in $`\delta _j(t)`$ terms, it becomes
$$\ddot{𝜹}=\left(\frac{𝑭}{𝑿}\right)|_{𝑿=𝑪}𝜹=\mathrm{J}\left[𝑪(t)\right]𝜹,$$
(27)
where $`\mathrm{J}\left[𝑪(t)\right]`$ is the Jacobi matrix. The latter system is also invariant under the action of the operator $`\widehat{g}=\widehat{a}`$, but its transformation must be correctly written as follows
$$\ddot{𝜹}=\widehat{g}^1\left(\frac{𝑭}{𝑿}\right)|_{𝑿=\widehat{g}𝑪}\widehat{g}𝜹=\widehat{g}^1\mathrm{J}\left[\widehat{g}𝑪(t)\right]\widehat{g}𝜹.$$
(28)
In other words, we have to replace the vector $`𝑿`$ in the Jacobi matrix by a *transformed* vector, $`\widehat{g}𝑪`$, near which the linearization is performed. Thus, we must write this matrix in the form $`\mathrm{J}\left[\widehat{g}𝑪(t)\right]`$ instead of $`\mathrm{J}\left[𝑪(t)\right]`$. In our case, $`\widehat{g}𝑪(t)\widehat{a}𝑪(t)=\{A(t),A(t),A(t),A(t)\}`$ and, therefore, we indeed have to add the above mentioned artificial transformation $`A(t)A(t)`$).
On the other hand, dealing with the linearized system $`\ddot{𝜹}=\mathrm{J}(t)𝜹`$, we conventionally consider the Jacobi matrix $`\mathrm{J}\left[𝑪(t)\right]\mathrm{J}(t)`$ as a fixed (but depending on $`t`$) matrix which *does not change* when the operator $`\widehat{g}=\widehat{a}`$ acts on the system $`\ddot{𝜹}=\mathrm{J}(t)𝜹`$ — this operator acts on the vector $`𝜹`$ only! The fact is that we try to split the system $`\ddot{𝜹}=\mathrm{J}(t)𝜹`$ into some subsystems using the *traditional algebraic transformations* of the old variables $`\delta _j`$. Indeed, we introduce new variables $`𝜹_{new}=\mathrm{S}𝜹`$, where $`\mathrm{S}`$ is a suitable time-independent orthogonal matrix, and then obtain the new system $`\ddot{𝜹}_{new}=\left(\stackrel{~}{\mathrm{S}}\mathrm{J}(t)\mathrm{S}\right)𝜹_{new}`$ that decomposes into a number of subsystems.
### II.4 Stability of the bush B$`[\widehat{a}^2,\widehat{ı}]`$ for the FPU-$`\alpha `$ chain with $`N>4`$ particles
Linearizing the dynamical equations of the FPU-$`\alpha `$ chain with $`N=6`$ in the vicinity of the bush B$`[\widehat{a}^2,\widehat{ı}]`$ ($`\pi `$-mode), we obtain the following Jacobi matrix in Eq. (20):
$$\mathrm{J}(t)=\mathrm{L}+4A(t)\mathrm{M},$$
where
$$\mathrm{L}=\left(\begin{array}{cccccc}\hfill 2& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 1\\ \hfill 1& \hfill 2& \hfill 1& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 1& \hfill 2& \hfill 1& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 1& \hfill 2& \hfill 1& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 2& \hfill 1\\ \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 2\end{array}\right),\mathrm{M}=\left(\begin{array}{cccccc}\hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 1\\ \hfill 1& \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 1& \hfill 0& \hfill 1& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 1& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 1\\ \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 0\end{array}\right).$$
These two symmetric matrices, unlike the case $`N=4`$, *do not* commute with each other:
$$\mathrm{L}\mathrm{M}\mathrm{M}\mathrm{L}=8A(t)\left(\begin{array}{cccccc}\hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 1& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 1\\ \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 0\\ \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 1\\ \hfill 1& \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 1& \hfill 0& \hfill 1& \hfill 0& \hfill 0\end{array}\right)0.$$
As a consequence, we cannot diagonalize both matrices $`\mathrm{L}`$ and $`\mathrm{M}`$ simultaneously, i. e. with the aid of one and the same orthogonal matrix $`\mathrm{S}`$. Therefore, it is impossible to diagonalize the Jacobi matrix $`\mathrm{J}(t)`$ in equation $`\ddot{𝜹}=\mathrm{J}(t)𝜹`$ for all time $`t`$. In other words, there are no such matrix $`\mathrm{S}`$ that completely splits the linearized system for the bush B$`[\widehat{a}^2,\widehat{ı}]`$ for the chain with $`N=6`$ particles.
This difference between the cases $`N=4`$ and $`N=6`$ (generally, for $`N>4`$) can be explained as follows. The group $`G=[\widehat{a}^2,\widehat{ı}]`$ of the considered bush, in fact, determines *different* groups for the cases $`N=4`$ and $`N=6`$. Indeed, for $`N=4`$ $`[\widehat{a}^2,\widehat{ı}]\{\widehat{E},\widehat{a}^2,\widehat{ı},\widehat{a}^2\widehat{ı}\}=D_2`$, while for $`N=6`$ $`[\widehat{a}^2,\widehat{ı}]\{\widehat{E},\widehat{a}^2,\widehat{a}^4,\widehat{ı},\widehat{a}^2\widehat{ı},\widehat{a}^4\widehat{ı}\}=D_3`$. The latter group ($`D_3`$) is non-Abelian ($`\widehat{ı}\widehat{a}^4=a^2\widehat{ı}`$), unlike the group $`D_2`$ ($`\widehat{ı}\widehat{a}^2=a^2\widehat{ı}`$) and, as a consequence, it possesses not only one-dimensional irreducible representations, but two-dimensional irreps, as well. It will be shown in Sec. IV, that precisely this fact does not permit us to split fully the above discussed linearized system <sup>3</sup><sup>3</sup>3Actually, this fact can be understood, if one takes into account that the two-dimensional irrep contains two times in the decomposition of the mechanical representation of the considered chain..
In spite of this difficulty, we can simplify the linearized system $`\ddot{𝜹}=\mathrm{J}(t)𝜹`$ considerably with the aid of some group-theoretical methods, which are discussed in the two following sections. Now, we only would like to present the final result of the above splitting for the case $`N=6`$:
$`\ddot{y}_1=4y_1,`$ (29a)
$`\ddot{y}_2=0,`$ (29b)
$`\{\begin{array}{c}\ddot{y}_3+2y_3=P(t)y_5,\hfill \\ \ddot{y}_5+2y_5=\overline{P}(t)y_3,\hfill \end{array}`$ (29e)
$`\{\begin{array}{c}\ddot{y}_4+2y_4=P(t)y_6,\hfill \\ \ddot{y}_6+2y_6=\overline{P}(t)y_4.\hfill \end{array}`$ (29h)
Here $`P(t)=e^{\frac{i\pi }{3}}4A(t)[1+e^{\frac{i\pi }{3}}]`$, while $`\overline{P}(t)`$ is the complex conjugate function with respect to $`P(t)`$. The two-dimensional subsystems (29e) and (29h) can be reduces to the real form (87) (see, Sec. V.1) by a certain linear transformation.
Note, that the stability of the $`\pi `$-mode (the bush B$`[\widehat{a}^2,\widehat{ı}]`$) was discussed in a number of papers \[Bud, ; Sand, ; Flach, , FPU1, ; FPU2, ; PR, ; Yosh, ; Shin, ; CLL, ; AntiFPU, \] by different methods and with an emphasis on different aspects of this stability. In particular, in our paper FPU1 , a remarkable fact was revealed for the FPU-$`\alpha `$ chain: the stability threshold of the $`\pi `$-mode is one and the same for interactions with *all* the other modes of the chain. (For other one-dimensional nonlinear modes, for both the FPU-$`\alpha `$ and FPU-$`\beta `$ chains, the stability thresholds, determined by interactions with different modes, are essentially different FPU2 ).
## III The general theorem and its consequence
We consider an $`N`$-degrees-of-freedom mechanical system that described by $`N`$ autonomous differential equations
$$\ddot{𝑿}=𝑭(𝑿),$$
(30)
where the configuration vector $`𝑿=\{x_1(t),x_2(t),\mathrm{},x_N(t)\}`$ determines the deviation from the equilibrium state $`𝑿=\{0,0,\mathrm{},0\}`$, while vector-function $`𝑭(𝑿)=\{f_1(𝑿),f_2(𝑿),\mathrm{},f_N(𝑿)\}`$ determines the right-hand-sides of the dynamical equations.
We assume that Eq. (30) is invariant under the action of a discrete symmetry group $`G_0`$ which we call “the parent symmetry group” of our mechanical system. This means that for all $`gG_0`$ Eq. (30) is invariant under the transformation of variables
$$\stackrel{~}{𝑿}=\widehat{g}𝑿,$$
(31)
where $`\widehat{g}`$ is the operator associated with the symmetry element $`g`$ of the group $`G_0`$ by the conventional definition
$$\widehat{g}𝑿=\{g^1x_1(t),\mathrm{},g^1x_N(t)\}.$$
Using (30) and (31), one can write $`𝑿=\widehat{g}^1\stackrel{~}{𝑿}`$, $`\widehat{g}^1\ddot{\stackrel{~}{𝑿}}=𝑭(\widehat{g}^1\stackrel{~}{𝑿})`$, and finally
$$\ddot{\stackrel{~}{𝑿}}=\widehat{g}𝑭(\widehat{g}^1\stackrel{~}{𝑿}).$$
(32)
On the other hand, renaming $`𝑿`$ from Eq. (30) as $`\stackrel{~}{𝑿}`$, one can write $`\ddot{\stackrel{~}{𝑿}}=𝑭(\stackrel{~}{𝑿})`$. Comparing this equation with Eq. (32), we obtain $`𝑭(\stackrel{~}{𝑿})=\widehat{g}𝑭(\widehat{g}^1\stackrel{~}{𝑿})`$, or
$$𝑭(\widehat{g}𝑿)=\widehat{g}𝑭(𝑿).$$
(33)
This is the condition of invariance of the dynamical equations (30) under the action of the operator $`\widehat{g}`$. It must hold for all $`gG_0`$ (obviously, it is sufficient to consider such equivalence only for the *generators* of the group $`G_0`$).
Let $`𝑿(t)=𝑪(t)`$ be an $`m`$-dimensional specific dynamical regime in the considered mechanical system that corresponds to the bush B$`[G]`$ ($`GG_0`$). This means that there exist some functional relations between the individual displacements $`x_i(t)`$ ($`i=1,2,\mathrm{},N`$), and, as a result, the system (30) reduces to $`m`$ ordinary differential equations in terms of the independent functions (we denoted them by $`A(t)`$, $`B(t)`$, $`C(t)`$, etc. in the previous section, see, for example, Eqs. (7,8)).
The vector $`𝑪(t)`$ is a general solution to the equation (see, Eq. (17))
$$\widehat{G}𝑿=𝑿,$$
where $`G`$ is the symmetry group of the given bush B$`[G]`$ ($`GG_0`$).
Now, we want to study the *stability* of the dynamical regime $`𝑪(t)`$, corresponding to the bush B$`[G]`$. To this end, we must linearize the dynamical equations (30) in a vicinity of the given bush, or more precisely, in a vicinity of the vector $`𝑪(t)`$. Let
$$𝑿=𝑪(t)+𝜹(t),$$
(34)
where $`𝜹(t)=\{\delta _1(t),\mathrm{},\delta _N(t)\}`$ is an infinitesimal $`N`$-dimensional vector. Substituting $`𝑿(t)`$ from (34) into (30) and linearizing these equations with respect to $`𝜹(t)`$, we obtain
$$\ddot{𝜹}=\mathrm{J}[𝑪(t)]𝜹,$$
(35)
where $`\mathrm{J}[𝑪(t)]`$ is the Jacobi matrix of the system (30):
$$\mathrm{J}[𝑪(t)]=\frac{f_i}{x_j}|_{𝑿=𝑪(t)}.$$
Now, we intend to prove the following
###### Theorem 1.
The matrix $`\mathrm{J}[𝐂(t)]`$ of the linearized dynamical equations near a given bush B$`[G]`$, determined by the configuration vector $`𝐂(t)`$, commutes with all matrices $`\mathrm{M}(g)`$ ($`gG`$) of the mechanical representation of the symmetry group $`G`$ of the considered bush:
$$\mathrm{M}(g)\mathrm{J}[𝑪(t)]=\mathrm{J}[𝑪(t)]\mathrm{M}(g).$$
###### Proof.
As was already discussed in Sec. II, the original nonlinear system $`\ddot{𝑿}=𝑭(𝑿)`$ transforms into the system $`\ddot{𝑿}=\widehat{g}^1𝑭(\widehat{g}𝑿)`$ under the action of the operator $`\widehat{g}`$ associated with the symmetry element $`gG_0`$ of the parent group $`G_0`$. According to Eq. (33), the invariance of our system with respect to the operator $`\widehat{g}`$ can be written as follows:
$$\widehat{g}^1𝑭(\widehat{g}𝑿)=𝑭(𝑿).$$
(36)
On the other hand, the system $`\ddot{𝑿}=𝑭(𝑿)`$, linearized in the vicinity of the vector $`𝑿=𝑪(t)`$ reads $`\ddot{𝜹}=\mathrm{J}\left[𝑪(t)\right]𝜹`$ (see Eq. (35)). Under the action of the operator $`\widehat{g}`$, it transforms, according to Eq. (28), into the system
$$\ddot{𝜹}=\widehat{g}^1\mathrm{J}\left[\widehat{g}𝑪(t)\right]\widehat{g}𝜹.$$
(37)
Let us now consider the mechanical representation $`\mathrm{\Gamma }`$ of the parent symmetry group $`G_0`$. To this end, we chose the “natural” basis $`𝚽=\{𝒆_1,𝒆_2,\mathrm{},𝒆_N\}`$ in the space of all possible displacements of individual particles (configuration space):
$$𝒆_1=\left(\begin{array}{c}1\\ 0\\ 0\\ \mathrm{}\\ 0\end{array}\right),𝒆_2=\left(\begin{array}{c}0\\ 1\\ 0\\ \mathrm{}\\ 0\end{array}\right),\mathrm{},𝒆_N=\left(\begin{array}{c}0\\ 0\\ 0\\ \mathrm{}\\ 1\end{array}\right).$$
(38)
Acting by an operator $`\widehat{g}`$ ($`gG`$) on the vector $`𝒆_j`$, we can write
$$\widehat{g}𝒆_j=\underset{i=1}{\overset{N}{}}\mathrm{M}_{ij}(g)𝒆_i,j=1,2,\mathrm{},N.$$
(39)
This equation associates the matrix $`\mathrm{M}(g)\mathrm{M}_{ij}`$ with the operator $`\widehat{g}`$ and, therefore, with the symmetry element $`gG`$:
$$g\widehat{g}\mathrm{M}(g).$$
(40)
The set of matrices $`\mathrm{M}(g)`$ corresponding to all $`gG`$ forms the mechanical representation $`\mathrm{\Gamma }`$ for our system <sup>4</sup><sup>4</sup>4According to the traditional definition of the $`n`$-dimensional matrix representation of the group $`G`$, a matrix $`\mathrm{M}(g)`$ is associated with the element $`gG`$, if $`\widehat{g}𝚽=\stackrel{~}{\mathrm{M}}(g)𝚽`$. Here $`𝚽=\{\mathit{\varphi }_1(𝒓),\mathit{\varphi }_2(𝒓),\mathrm{},\mathit{\varphi }_N(𝒓)\}`$ is the set of basis vectors, $`\widehat{g}`$ is the operator acting on the vectors as $`\widehat{g}\mathit{\varphi }_i(𝒓)=\mathit{\varphi }_i(g^1𝒓)`$, and $`\stackrel{~}{\mathrm{M}}(g)`$ is the matrix transposed with respect to the matrix $`\mathrm{M}(g)`$.. As a consequence of this definition, the equation
$$\widehat{g}𝑪=\mathrm{M}(g)𝑪$$
(41)
is valid for any vector $`𝑪`$ determined in the basis (38) as $`𝑪=_{k=1}^NC_k𝒆_k`$.
Using Eq. (41), we can rewrite the equation (37) in terms of matrices $`\mathrm{M}(g)\mathrm{M}_g`$ ($`gG_0`$) of the mechanical representation of the group $`G_0`$:
$$\ddot{𝜹}=\mathrm{M}_g^1\mathrm{J}\left[\mathrm{M}_g𝑪(t)\right]\mathrm{M}_g𝜹.$$
(42)
Therefore, the invariance of the system $`\ddot{𝜹}=\mathrm{J}\left[𝑪(t)\right]𝜹`$ with respect of the operator $`\widehat{g}`$ (matrix $`\mathrm{M}_g`$) can be written as the following relation
$$\mathrm{M}_g^1\mathrm{J}\left[\mathrm{M}_g𝑪(t)\right]\mathrm{M}_g=\mathrm{J}\left[𝑪(t)\right].$$
(43)
Now, let us suppose that $`g`$ is an element of the symmetry group $`G`$ of a given bush B$`[G]`$ ($`GG_0`$). By the definition, all the elements of this group ($`gG`$) leave invariant the vector $`𝑪(t)`$ that determines the displacement pattern of this bush
$$\widehat{g}𝑪(t)=\mathrm{M}_g𝑪(t)=𝑪(t),gG.$$
(44)
Taking into account this equation, we obtain from (43) the relation
$$\mathrm{M}_g^1\mathrm{J}\left[𝑪(t)\right]\mathrm{M}_g=\mathrm{J}\left[𝑪(t)\right],$$
(45)
which holds for each element $`g`$ of the symmetry group $`G`$ of the considered bush.
Rewriting (45) in the form
$$\mathrm{J}\left[𝑪(t)\right]\mathrm{M}_g=\mathrm{M}_g\mathrm{J}\left[𝑪(t)\right],$$
(46)
we arrive at the conclusion of our Theorem: all the matrices $`\mathrm{M}_g`$ of the mechanical representation of the group $`G`$ *commute* with the Jacobi matrix $`\mathrm{J}\left[𝑪(t)\right]`$ of the linearized (near the given bush) dynamical equations $`\ddot{𝜹}=\mathrm{J}\left[𝑪(t)\right]𝜹`$. ∎
In what follows, we will introduce a simpler notation for the Jacobi matrix:
$$\mathrm{J}\left[𝑪(t)\right]\mathrm{J}(t).$$
(47)
*Remark*. We have proved that all the matrices $`\mathrm{M}_g`$ with $`gG`$ commute with the Jacobi matrix $`\mathrm{J}(t)`$ of the system (35). But if we take a symmetry element $`gG_0`$ that is not contained in $`G`$ ($`gG_0G`$), the matrix $`\mathrm{M}_g`$ corresponding to $`g`$ *may not* commute with $`\mathrm{J}(t)`$. An example of such noncommutativity and the source of this phenomenon were presented in Sec. II.3.
#### Consequence of Theorem 1
Taking into account Theorem 1, we can apply the well-known *Wigner* theorem Dob to split the linearized system $`\ddot{𝜹}=\mathrm{J}(t)𝜹`$ into a certain number of independent subsystems. Indeed, according to this theorem, the matrix ($`\mathrm{J}(t)`$, in our case) commuting with all the matrices of a representation $`\mathrm{\Gamma }`$ of the group $`G`$ (mechanical representation, in our case), can be reduced to a very specific block-diagonal form. The dimension of each block of this form is equal to $`n_jm_j`$, where $`n_j`$ is the dimension of a certain *irreducible representation* (irrep) $`\mathrm{\Gamma }_j`$ of the group $`G`$ containing $`m_j`$ times in the reducible representation $`\mathrm{\Gamma }`$. Moreover, these blocks possess a particular structure which will be considered in Sec. IV.
To implement this splitting explicitly one must pass from the old basis $`𝚽_{old}=\{𝒆_1,𝒆_2,\mathrm{},𝒆_N,\}`$ of the mechanical space to the new basis $`𝚽_{new}=\{\mathit{\varphi }_1,\mathit{\varphi }_2,\mathrm{},\mathit{\varphi }_N,\}`$ formed by the complete set of the basis vectors $`\mathit{\varphi }_k`$ ($`k=1,2,\mathrm{},N`$) of all the irreps of the group $`G`$. If $`𝚽_{new}=\mathrm{S}𝚽_{old}`$ then the unitary transformation <sup>5</sup><sup>5</sup>5Here $`\mathrm{S}^+`$ is the Hermite conjugated matrix with respect to the matrix $`\mathrm{S}`$
$$\mathrm{J}_{new}(t)=\mathrm{S}^+\mathrm{J}_{old}(t)\mathrm{S}$$
(48)
produces the above discussed block-diagonal matrix $`\mathrm{J}_{new}(t)`$ of the linearized system $`\ddot{𝜹}_{new}=\mathrm{J}_{new}(t)𝜹_{new}`$ (here $`𝜹_{old}=\mathrm{S}𝜹_{new}`$).
In the next section, we will search the basis vectors $`\mathit{\varphi }_i[\mathrm{\Gamma }_j]`$ ($`i=1,2,\mathrm{},n_j`$) of each irreducible representation $`\mathrm{\Gamma }_j`$ in the form
$$\mathit{\varphi }_i[\mathrm{\Gamma }_j]=\{x_{ij}^{(1)},x_{ij}^{(2)},\mathrm{},x_{ij}^{(N)}\},$$
(49)
where $`x_{ij}^{(k)},`$ ($`k=1,2,\mathrm{},N`$) determines the displacement of the $`k`$-th particle corresponding to the $`i`$-th basis vector $`\mathit{\varphi }_i[\mathrm{\Gamma }_j]`$ of the $`j`$-th irrep $`\mathrm{\Gamma }_j`$. Actually, this means that we search $`\mathit{\varphi }_i[\mathrm{\Gamma }_j]`$ as a superposition of the old basis vectors $`𝒆_k`$ ($`k=1,2,\mathrm{},N`$) of the mechanical space (see (38)):
$$\mathit{\varphi }_i[\mathrm{\Gamma }_j]=\underset{k=1}{\overset{N}{}}x_{ij}^{(k)}𝒆_k.$$
(50)
If we find all the basis vectors $`\mathit{\varphi }_i[\mathrm{\Gamma }_j]`$ in such a form, the coefficients $`x_{ij}^{(k)}`$ are obviously the elements of the matrix $`\mathrm{S}`$ that determines the transformation $`𝚽_{new}=\mathrm{S}𝚽_{old}`$ from the old basis $`𝚽_{old}=\{𝒆_k|k=1,2,\mathrm{},N\}`$ to the new basis $`𝚽_{new}=\{\mathit{\varphi }_i[\mathrm{\Gamma }_j]|i=1,2,\mathrm{},n_j;j=1,2,\mathrm{}\}`$. Here $`j=1,2,\mathrm{}`$ are indices of the irreducible representations that contribute to the reducible mechanical representation $`\mathrm{\Gamma }`$.
Thus, finding all the basis vectors $`\mathit{\varphi }_i[\mathrm{\Gamma }_j]`$ of the irreps $`\mathrm{\Gamma }_j`$ in the form (49) provides us directly with the matrix $`\mathrm{S}`$ that diagonalizes the Jacobi matrix $`\mathrm{J}(t)`$ of the linearized system $`\ddot{𝜹}=\mathrm{J}(t)𝜹`$.
## IV Stability analysis of dynamical regimes in monoatomic chains
### IV.1 Setting up the problem and Theorem 2
In general, the study of stability of periodic and, especially, quasiperiodic dynamical regimes in the mechanical systems with many degrees of freedom presents considerable difficulties. Indeed, for this purpose, we must integrate large linearized (near the considered regime) system of differential equations with time-dependent coefficients. In the case of periodic regime, one can use the Floquet method requiring integration over only one time-period to construct the monodromy matrix. But for quasiperiodic regime this method is inapplicable, and one often needs to solve system of great number of differential equations for very large time intervals to reveal instability (especially, near the stability threshold).
In such a situation, a decomposition (splitting) of the full linearized system into a number of independent subsystems of small dimensions proves to be very useful. Moreover, this decomposition can provide valuable information on generalized degrees of freedom responsible for the loss of stability of the given dynamical regime for the first time. Let us note that the number of such “critical” degrees of freedom can frequently be rather small.
We want to illustrate the above idea with the case of $`N`$-particle monoatomic chains for $`N1`$. Let us introduce the following notation. The bush B$`[G]`$ with the symmetry group $`G`$ containing the translational subgroup $`[\widehat{a}^m]`$ will be denoted by B$`[\widehat{a}^m,\mathrm{}]`$, where dots stand for other generators of the group $`G`$. (Note, that any $`m`$-dimensional bush can exist only for the chain with $`N`$ divisible by $`m`$).
###### Theorem 2.
Linear stability analysis of any bush B$`[\widehat{a}^m,\mathrm{}]`$ in the $`N`$-degrees-of-freedom monoatomic chain can be reduced to stability analysis of isolated subsystems of the second order differential equations with time-dependent coefficients whose dimensions do not exceed the integer number $`m`$.
*Corollary*. If the bush dimension is $`d`$, one can pass on to the subsystems of *autonomous* differential equations with dimensions not exceeding $`(m+d)`$.
Before proving these propositions we must consider the procedure of constructing the basis vectors of the irreducible representations of the translational group $`T`$.
### IV.2 Basis vectors of irreducible representations of the translational groups
The basis vectors of irreducible representations of different symmetry groups are usually obtained by the method of projection operators Dob , but, in our case, it is easier to make use of the “direct” method based on the definition of the group representation <sup>6</sup><sup>6</sup>6We already use this method in our previous papers (see, for example, FPU1 ).
Let $`\mathrm{\Gamma }`$ be an $`n`$-dimensional representation (reducible or irreducible) of the group $`G`$, while $`V[\mathrm{\Gamma }]`$ be the invariant subspace corresponding to this representation that determined by the set $`𝚽`$ of $`N`$-dimensional basis vectors $`\mathit{\varphi }_j`$ ($`j=1,\mathrm{},n`$):
$$𝚽=\{\mathit{\varphi }_1,\mathit{\varphi }_2,\mathrm{},\mathit{\varphi }_n\}.$$
(51)
Acting on any basis vector $`\mathit{\varphi }_j`$ by an operator $`\widehat{g}`$ ($`gG`$) and bearing in mind the invariance of the subspace $`V[\mathrm{\Gamma }]`$, we can represent the vector $`\widehat{g}\mathit{\varphi }_j`$ as a superposition of all basis vectors from (51). In other words,
$$\widehat{g}𝚽\{\widehat{g}\mathit{\varphi }_1,\widehat{g}\mathit{\varphi }_2,\mathrm{},\widehat{g}\mathit{\varphi }_n\}=\stackrel{~}{\mathrm{M}}(g)𝚽,$$
(52)
where $`\mathrm{M}(g)`$ is the matrix corresponding, in the representation $`\mathrm{\Gamma }`$, to the element $`g`$ of the group $`G`$. (In Eq. (52) we use tilde as the symbol of matrix transposition). Eq. (52) associates with any $`gG`$ a certain $`n\times n`$ matrix $`\mathrm{M}(g)`$ and encapsulates the definition of matrix representation
$$\mathrm{\Gamma }=\{\mathrm{M}(g_1),\mathrm{M}(g_2),\mathrm{}\}.$$
(53)
The above mentioned “direct” method is based precisely on this definition. Let us use it to obtain the basis vectors of the irreducible representations for the translational group $`T[\widehat{a}^m]`$. We will construct these vectors in the mechanical space of the $`N`$-particle monoatomic chain and, therefore, each vector $`\mathit{\varphi }_j`$ can be written as follows:
$$\mathit{\varphi }_j=\{x_1,x_2,\mathrm{},x_N\},$$
(54)
where $`x_i`$ is a displacement of $`i`$-th particle from its equilibrium.
The group $`T[\widehat{a}^m]`$ represents a *translational subgroup* corresponding to the bush B$`[G]=`$B$`[\widehat{a}^m,\mathrm{}]`$. For $`N`$-particle chain <sup>7</sup><sup>7</sup>7Note, the relation $`Nmodm=0`$ must hold!, $`T[\widehat{a}^m]`$ is a subgroup of the order $`k=N/m`$ of the full translational group $`T_N[\widehat{a}]`$, and we can write the complete set of its elements as follows:
$$T_k=\{\widehat{e},\widehat{a}^m,\widehat{a}^{2m},\widehat{a}^{3m},\mathrm{},\widehat{a}^{(k1)m}\}(\widehat{a}^{km}=\widehat{a}^N=\widehat{e}).$$
(55)
Being cyclic, the group $`T_k`$ from (55) possesses only one-dimensional irreps, and their total number is equal to the order ($`k=N/m`$) of this group.
Below, for simplicity, we consider the case $`m=3`$ and $`N=12`$. The generalization to the case of arbitrary values of $`m`$ and $`N`$ turns out to be trivial.
As it is well-known, the one-dimensional irreps $`\mathrm{\Gamma }_i`$ of the $`k`$-order cyclic group can be constructed with the aid of $`k`$-degree roots of $`1`$ and, therefore, for our case $`N=12`$, $`m=3`$, $`k=4`$, we obtain the irreducible representations listed in Table 1.
In accordance with the definition (52), the basis vector $`\mathit{\varphi }`$ of the one-dimensional irrep $`\mathrm{\Gamma }`$, for which $`\mathrm{M}(g)=\gamma `$, must satisfy the equation
$$\widehat{g}\mathit{\varphi }=\gamma \mathit{\varphi }.$$
(56)
In our case, $`\widehat{g}=\widehat{a}^3`$, this equation can be written as follows:
$$\begin{array}{c}\widehat{g}\mathit{\varphi }\{x_{10},x_{11},x_{12}|x_1,x_2,x_3|x_4,x_5,x_6|x_7,x_8,x_9\}=\hfill \\ \gamma \{x_1,x_2,x_3|x_4,x_5,x_6|x_7,x_8,x_9|x_{10},x_{11},x_{12}\}.\hfill \end{array}$$
(57)
Here $`\gamma =1,i,1,i`$ for the irreps $`\mathrm{\Gamma }_1`$, $`\mathrm{\Gamma }_2`$, $`\mathrm{\Gamma }_3`$ and $`\mathrm{\Gamma }_4`$, respectively. Equating the sequential components of both sides of Eq. (57), we obtain the *general solution* to the equation $`\widehat{g}\mathit{\varphi }=\gamma \mathit{\varphi }`$ that turns out to depend on three arbitrary constants, say, $`x`$, $`y`$ and $`z`$:
$$\begin{array}{c}\mathit{\varphi }=\{x,y,z|\gamma ^1x,\gamma ^1y,\gamma ^1z|\gamma ^2x,\gamma ^2y,\gamma ^2z|\gamma ^3x,\gamma ^3y,\gamma ^3z\}=\hfill \\ x\{1,0,0|\gamma ^1,0,0|\gamma ^2,0,0|\gamma ^3,0,0\}+\hfill \\ y\{0,1,0|0,\gamma ^1,0|0,\gamma ^2,0|0,\gamma ^3,0\}+\hfill \\ z\{0,0,1|0,0,\gamma ^1|0,0,\gamma ^2|0,0,\gamma ^3\}.\hfill \end{array}$$
(58)
Here we write the vector $`\mathit{\varphi }`$ as the superposition (with coefficients $`x`$, $`y`$, $`z`$) of *three basis vectors*. It means that the irrep $`\mathrm{\Gamma }`$ is contained *thrice* in the decomposition of the mechanical representation into irreducible representations of the group $`G=[\widehat{a}^3]`$.
This result can be generalized to the case of arbitrary $`N`$ and $`m`$ in trivial manner: each irrep of the group $`G=[\widehat{a}^m]`$ enters *exactly* $`m`$ *times* into the decomposition of the mechanical representation for $`N`$-particle chain, and the rule for constructing $`m`$ appropriate basis vectors is fully obvious from Eq. (58).
### IV.3 Proof of Theorem 2
###### Proof.
The basis vectors of all irreps $`\mathrm{\Gamma }_i`$, listed for the case $`N=12`$, $`m=3`$ in Table 1, can be obtained from (58) setting $`\gamma =1,i,1,i`$, respectively (these values are one-dimensional matrices corresponding in $`\mathrm{\Gamma }_i`$ ($`i=1,2,3,4`$) to the generator $`\widehat{g}\widehat{a}^3`$).
Let us write the above basis vectors sequentially, as it is done in Table 2, and prove that $`12\times 12`$ matrix, determined by this table, is precisely the matrix $`\mathrm{S}`$ that splits the linearized dynamical equations $`\ddot{𝜹}=\mathrm{J}(t)𝜹`$ for the considered case. In Table 2, we denote the basis vectors $`\mathit{\varphi }_j(\mathrm{\Gamma }_i)`$ by the symbol of the irrep $`\mathrm{\Gamma }_i`$ ($`i=1,2,3,4`$) and the number $`j=1,2,3`$ of the basis vector of this irrep. The normalization factor $`\left(\frac{1}{2}\right)`$ must be associated with each row of this table to produce the normalized basis vectors (because of this fact, we mark the rows as $`2\mathit{\varphi }_j(\mathrm{\Gamma }_i)`$ in the last column of Table 2).
Obviously, we can use the matrix $`\mathrm{S}`$ from Table 2 (the rows of this matrix are the basis vectors of all the irreps of the group $`T_4`$) not only for the action on the vectors in the $`𝑿`$-space of the full nonlinear system, but on the vectors in the $`𝜹`$-space of the linearized system $`\ddot{𝜹}=\mathrm{J}(t)𝜹`$, as well. It is essential that in the latter case the matrix $`\mathrm{S}`$ reduces the Jacobi matrix $`\mathrm{J}(t)`$ to a certain block-diagonal form. Indeed, as was shown in Theorem 1, the matrix $`\mathrm{J}(t)`$ commutes with all the matrices of the mechanical representation of the bush symmetry group. Therefore, according to Wigner theorem Dob , it can be reduced, using unitary transformation by the matrix $`\mathrm{S}`$, to the block-diagonal form with blocks whose dimension is equal to $`n_jm_j`$. Here $`n_j`$ is the dimension of the irrep $`\mathrm{\Gamma }_j`$, while $`m_j`$ is the number of times that this irrep enters into the decomposition of the mechanical representation (constructed, in our case, in the $`𝜹`$-space).
In Sec. IV.2, we have shown that for the translational group $`T=[\widehat{a}^m]`$ all $`n_j=1`$ and all $`m_j=m`$. Therefore, the above matrix $`\mathrm{S}`$ decomposes the Jacobi matrix $`\mathrm{J}(t)`$ into blocks whose dimension is equal to $`m`$. As a consequence of this decomposition, the system $`\ddot{𝜹}=\mathrm{J}(t)𝜹`$ splits into $`k=N/m`$ independent subsystems, $`L_j`$ ($`j=1,2,\mathrm{},k`$), each consisting of $`m`$ differential equations of the second order. The coefficients of these equations are time-dependent functions, and this time dependence is determined by the functions $`A(t)`$, $`B(t)`$, $`C(t)`$ etc., entering into the bush displacement pattern (see, for example, (7,8)). For the one-dimensional bushes, the coefficients of the above subsystems $`L_j`$ turn out to be periodic functions with identical period, while for the many-dimensional bushes they possess different periods (such bushes describe quasiperiodic motion).
In general, it is impossible to obtain the explicit form of the functions $`A(t)`$, $`B(t)`$, $`C(t)`$ etc., determining the bush displacement pattern. Therefore, we can add the bush dynamical equations to the $`m`$ differential equations of each subsystems $`L_j`$. These additional $`d`$ equations determine the functions $`A(t)`$, $`B(t)`$, $`C(t)`$ etc. *implicitly*, where $`d`$ is the dimension of the considered bush B$`[\widehat{a}^m,\mathrm{}]`$.
On the other hand, we can give the following estimate for the bush dimension $`d`$:
$$dm.$$
(59)
Here $`m`$ is the index of the translational symmetry of the bush B$`[\widehat{a}^m,\mathrm{}]`$ (it determines the ratio between the size of the primitive cell in the vibrational state and in the equilibrium). Indeed, the bush displacement pattern can be found as the solution to the equation $`\widehat{G}𝑿=𝑿`$. If we take into account only translational symmetry group of the bush B$`[\widehat{a}^m,\mathrm{}]`$, i.e. $`G=[\widehat{a}^m]`$, this equation reduces to the equation $`\widehat{g}\mathit{\varphi }=\mathit{\varphi }`$ ($`\widehat{g}=\widehat{a}^m`$) for the basis vector $`\mathit{\varphi }`$ of the *identity* irrep ($`\gamma =1`$ in (56)) of the group $`T=[\widehat{a}^m]`$. As it has been already shown in Sec. IV.2, such vector $`\mathit{\varphi }`$ depends on exactly $`m`$ arbitrary parameters. But some additional symmetry elements, denoted by dots in the bush symbol B$`[\widehat{a}^m,\mathrm{}]`$, can lead to a decrease in the number $`m`$ of the above parameters. For example, the bush B$`[\widehat{a}^3,\widehat{ı}]`$ turns out to be one-dimensional, i.e., in this case, $`d=1`$, while $`m=3`$. Even the *vibrational* bush B$`[\widehat{a}^3]`$ turns out to be two-dimensional ($`d=2`$), if the condition of immobility of the mass center is taken into account. Thus, for all cases, $`dm`$, and we can state that $`m`$ equations of each $`L_j`$, extended by $`d`$ additional equations of the given bush, provide us with $`k=N/m`$ independent subsystems $`\stackrel{~}{L}_j`$ of ($`m+d`$) *autonomous* differential equations.
Taking into account the additional bush symmetry elements, denoted by dots in the symbol B$`[\widehat{a}^m,\mathrm{}]`$, leads not only to reducing the bush dimension, but to a further splitting of the above discussed subsystems $`L_j`$ (we consider this point in the next section). Thus, the linear stability analysis of the bush B$`[\widehat{a}^m,\mathrm{}]`$ in the $`N`$-particle chain indeed reduces to studying stability of individual subsystems whose dimension does not exceed $`(m+d)`$. This is the conclusion of Theorem 2 and, thus, we have completed the proof. ∎
### IV.4 Example 1: splitting the linearized system for the bush B$`[\widehat{a}^3]`$
We consider the splitting of the linearized system $`\ddot{𝜹}=\mathrm{J}(t)𝜹`$ for the bush B$`[\widehat{a}^3]`$ in a chain with $`N=12`$ particles. The original nonlinear system, for this case, reads
$$\begin{array}{c}\ddot{x}_i=f(x_{i+1}x_i)f(x_ix_{i1}),\hfill \\ i=1,2,\mathrm{},12(x_0=x_{12},x_{13}=x_1).\hfill \end{array}$$
(60)
The displacement pattern of the bush B$`[\widehat{a}^3]`$, obtained from the equation $`\widehat{a}^3𝑿=𝑿`$ reads:
$$𝑿=\{x(t),y(t),z(t)|x(t),y(t),z(t)|x(t),y(t),z(t)|x(t),y(t),z(t)\}.$$
(61)
Substituting this form of vibrational pattern into (60), we obtain three differential equations for the functions $`x(t)`$, $`y(t)`$, $`z(t)`$ (all the other equations of (60) turn out to be equivalent to these equations):
$$\begin{array}{c}\ddot{x}=f(yx)f(xz),\hfill \\ \ddot{y}=f(zy)f(yx),\hfill \\ \ddot{z}=f(xz)f(zy).\hfill \end{array}$$
(62)
The linearization of the Eqs. (60) near the dynamical regime determined by (61) leads to the system
$$\ddot{𝜹}=\mathrm{J}(t)𝜹$$
(63)
with the following Jacobi matrix:
$$\mathrm{J}(t)=\left(\begin{array}{cccccccccccc}\alpha & A& 0& 0& 0& 0& 0& 0& 0& 0& 0& B\\ A& \beta & C& 0& 0& 0& 0& 0& 0& 0& 0& 0\\ 0& C& \gamma & B& 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& B& \alpha & A& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& A& \beta & C& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& C& \gamma & B& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& B& \alpha & A& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& A& \beta & C& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& C& \gamma & B& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& B& \alpha & A& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& 0& A& \beta & C\\ B& 0& 0& 0& 0& 0& 0& 0& 0& 0& C& \gamma \end{array}\right),$$
(64)
where
$$\begin{array}{c}A(t)=f^{}[y(t)x(t)],\hfill \\ B(t)=f^{}[x(t)z(t)],\hfill \\ C(t)=f^{}[z(t)y(t)],\hfill \\ \alpha (t)=[A(t)+B(t)],\hfill \\ \beta (t)=[A(t)+C(t)],\hfill \\ \gamma (t)=[B(t)+C(t)].\hfill \end{array}$$
(65)
Using Table 2, the matrix $`\mathrm{S}`$ that splits up the system (63) can be written as follows:
$$\mathrm{S}=\frac{1}{2}\left(\begin{array}{cccc}\hfill \mathrm{I}& \hfill \mathrm{I}& \hfill \mathrm{I}& \hfill \mathrm{I}\\ \hfill \mathrm{I}& \hfill i\mathrm{I}& \hfill \mathrm{I}& \hfill i\mathrm{I}\\ \hfill \mathrm{I}& \hfill \mathrm{I}& \hfill \mathrm{I}& \hfill \mathrm{I}\\ \hfill \mathrm{I}& \hfill i\mathrm{I}& \hfill \mathrm{I}& \hfill i\mathrm{I}\end{array}\right),$$
(66)
where $`\mathrm{I}`$ is the $`3\times 3`$ identity matrix
$$\mathrm{I}=\left(\begin{array}{ccc}\hfill 1& \hfill 0& \hfill 0\\ \hfill 0& \hfill 1& \hfill 0\\ \hfill 0& \hfill 0& \hfill 1\end{array}\right).$$
With the aid of the unitary transformation
$$\mathrm{J}_{new}(t)=\mathrm{S}^+\mathrm{J}(t)\mathrm{S},$$
(67)
we obtain
$$\mathrm{J}_{new}(t)=\left(\begin{array}{cccc}\mathrm{D}_1& 0& 0& 0\\ 0& \mathrm{D}_2& 0& 0\\ 0& 0& \mathrm{D}_3& 0\\ 0& 0& 0& \mathrm{D}_4\end{array}\right),$$
(68)
where
$$\mathrm{D}_k=\left(\begin{array}{ccc}(A+B)& A& \gamma _kB\\ A& (A+C)& C\\ \overline{\gamma }_kB& C& (B+C)\end{array}\right),(k=1,2,3,4),$$
(69)
with $`\gamma _1=1`$, $`\gamma _2=i`$, $`\gamma _3=1`$, $`\gamma _4=i`$ ($`\overline{\gamma }_k`$ is the complex conjugate value of $`\gamma _k`$).
This means that the linear transformation
$$𝜹=\mathrm{S}𝜹_{new}$$
(70)
reduces the old equations (63) to the following form
$$\ddot{𝜹}_{new}=\mathrm{J}_{new}(t)𝜹_{new}$$
(71)
with block-diagonal matrix $`\mathrm{J}_{new}(t)=\mathrm{S}^+\mathrm{J}(t)\mathrm{S}`$ determined by Eqs. (68,69).
Assuming
$$𝜹_{new}=\{\delta _1^{(1)},\delta _2^{(1)},\delta _3^{(1)}|\delta _1^{(2)},\delta _2^{(2)},\delta _3^{(2)}|\delta _1^{(3)},\delta _2^{(3)},\delta _3^{(3)}|\delta _1^{(4)},\delta _2^{(4)},\delta _3^{(4)}\},$$
we can present (71) in a more explicit form:
$$\begin{array}{c}\ddot{\delta }_1^{(k)}=(A+B)\delta _1^{(k)}+A\delta _2^{(k)}+\gamma _kB\delta _3^{(k)},\hfill \\ \ddot{\delta }_2^{(k)}=A\delta _1^{(k)}(A+C)\delta _2^{(k)}+C\delta _3^{(k)},\hfill \\ \ddot{\delta }_3^{(k)}=\overline{\gamma }_kB\delta _1^{(k)}+C\delta _2^{(k)}(B+C)\delta _3^{(k)},\hfill \end{array}$$
(72)
where $`\gamma _k=1,i,1,i`$ for $`k=1,2,3,4`$, respectively.
Thus, we obtain four independent $`3\times 3`$ systems of linear differential equations with time-dependent coefficients $`A(t)`$, $`B(t)`$ and $`C(t)`$, which are determined by Eqs. (65).
Let us write down these equations for the FPU-$`\alpha `$ chain. For this case, the function $`f(x)`$ in (60) reads $`f(x)=x+x^2`$. Therefore, $`f^{}(x)=1+2x`$, and we obtain from (65)
$$\begin{array}{c}A(t)=1+2[y(t)x(t)],\hfill \\ B(t)=1+2[x(t)z(t)],\hfill \\ C(t)=1+2[z(t)y(t)].\hfill \end{array}$$
Substituting these functions into (72) one can finally obtain the following equations for the FPU-$`\alpha `$ chain:
$$\ddot{𝜹}^{(k)}=\mathrm{J}_k(t)𝜹^{(k)},$$
where
$$\mathrm{J}_k(t)=\left(\begin{array}{ccc}2(1+yz)& 1+2(yx)& \gamma _k[1+2(xz)]\\ 1+2(yx)& 2(1+zx)& 1+2(zy)\\ \overline{\gamma }_k[1+2(xz)]& [1+2(zy)]& 2(1+xy)\end{array}\right),(k=1,2,3,4).$$
(73)
Here $`\gamma _1=1`$, $`\gamma _2=i`$, $`\gamma _3=1`$, $`\gamma _4=i`$ and $`xx(t)`$, $`yy(t)`$, $`zz(t)`$ are functions determined by the dynamical equations of the bush B$`[\widehat{a}^3]`$ :
$$\begin{array}{c}\ddot{x}=(y2x+z)(1+yz),\hfill \\ \ddot{y}=(z2y+x)(1+zx),\hfill \\ \ddot{z}=(x2z+y)(1+xy).\hfill \end{array}$$
(74)
These equations can be obtained from (62) taking into account the relation $`f(x)=x+x^2`$ for the case of the FPU-$`\alpha `$ model.
*Remark*. According to (61), (62) (see, also (74)), the vibrational bush B$`[\widehat{a}^3]`$ is three-dimensional. However, it actually turns out to be a *two-dimensional* bush. Indeed, there is no *onsite* potential in the FPU-$`\alpha `$ chains and, therefore, the conservation law of the total momentum of such system holds. Assuming that the center of masses is fixed, we obtain an additional relation $`x(t)+y(t)+z(t)=0`$ which reduces the dimension of the bush B$`[\widehat{a}^3]`$ from $`3`$ to $`2`$.
### IV.5 Further decomposition of linearized systems based on higher symmetry groups
Up to this point, we have discussed the decomposition of the linearized system $`\ddot{𝜹}=\mathrm{J}(t)𝜹`$ using only the *translational* part of the bush symmetry group. In general, one can arrive at a more detailed splitting, if one takes into account the additional bush symmetries.
#### IV.5.1 Example 2: Splitting of the linearized system for the bush B$`[\widehat{a}^3,\widehat{ı}]`$
Let us consider the decomposition of the linearized system for the bush B$`[\widehat{a}^3,\widehat{ı}]`$ in the case of an arbitrary monoatomic chain. Since translational part of the symmetry group $`G=[\widehat{a}^3,\widehat{ı}]`$, which turns out to be the dihedral group, is the same as that of the early considered bush B$`[\widehat{a}^3]`$, we can take advantage of all the results obtained in Sec. IV.4 and add only some restrictions originating from the presence of the additional generator $`\widehat{ı}`$ of the group $`[\widehat{a}^3,\widehat{ı}]`$.
Substituting the vector $`𝑿(t)`$ in the form (61) into the equation $`\widehat{ı}𝑿(t)=𝑿(t)`$, we obtain $`z(t)x(t)`$, $`y(t)y(t)`$ and, therefore, $`y(t)0`$. The displacement pattern for the bush B$`[\widehat{a}^3,\widehat{ı}]`$ then can be written as follows:
$$𝑿(t)=\{x(t),0,x(t)|x(t),0,x(t)|x(t),0,x(t)|x(t),0,x(t)\}.$$
(75)
Thus, the bush B$`[\widehat{a}^3,\widehat{ı}]`$ turns out to be *one-dimensional*.
As a result of the substitution $`z(t)x(t)`$, $`y(t)0`$, three equations (62) reduce to only one equation
$$\ddot{x}=f(x)f(2x).$$
(76)
For the FPU-$`\alpha `$ chain (see Eqs. (74)), this equation transforms to
$$\ddot{x}+3x+3x^2=0.$$
(77)
Unlike the purely translational group $`[\widehat{a}^3]`$, of the three-dimensional bush B$`[\widehat{a}^3]`$, the symmetry group $`[\widehat{a}^3,\widehat{ı}]`$ of the one-dimensional bush B$`[\widehat{a}^3,\widehat{ı}]`$ is the dihedral group with another sets of the irreps and basis vectors. It can be shown that taking into account that $`[\widehat{a}^3,\widehat{ı}]`$ is the supergroup with respect to group $`[\widehat{a}^3]`$, allows one to obtain the following splitting scheme of the linearized system $`\ddot{𝜹}=\mathrm{J}(t)𝜹`$
$$\begin{array}{ccc}1:(\delta _1);\hfill & 1:(\delta _6);\hfill & 2:(\delta _2,\delta _3);\hfill \\ 2:(\delta _4,\delta _5);\hfill & 3:(\delta _7,\delta _9,\delta _{11});\hfill & 3:(\delta _8,\delta _{10},\delta _{12}).\hfill \end{array}$$
(78)
Here we present the dimension of each independent subsystem (before the colon) and the list of its variables (after the colon). From the scheme (78), one can see that two of four three-dimensional subsystems corresponding to the splitting provided by group $`[\widehat{a}^3]`$ (see Eqs. (72)) in the case of the supergroup $`[\widehat{a}^3,\widehat{ı}]`$ are decomposed into new independent subsystems of dimensions equal to $`1`$ and $`2`$. Below we explain how one can obtain the splitting schemes analogous to (78).
#### IV.5.2 General case
Now we consider the application of the Wigner theorem in case of an arbitrary bush B$`[G]`$.
Let us consider a matrix $`\mathrm{H}`$ commuting with all matrices of a reducible representation $`\mathrm{\Gamma }`$ of the group $`G`$ that can be decomposed into the irreps $`\mathrm{\Gamma }_j`$ of this group as follows
$$\mathrm{\Gamma }=\underset{j}{}^{{\scriptscriptstyle }}m_j\mathrm{\Gamma }_j.$$
(79)
According to the Wigner theorem, the matrix $`\mathrm{H}`$ can be reduced to a block-diagonal form with the blocks $`\mathrm{D}_j`$ of dimensions $`n_jm_j`$, corresponding to the each $`\mathrm{\Gamma }_j`$, with $`n_j`$ being the dimension of the irrep $`\mathrm{\Gamma }_j`$ entering $`m_j`$ times into the decomposition (79) of the representation $`\mathrm{\Gamma }`$ into the irreducible parts.
Moreover, each block $`\mathrm{D}_j`$ possesses a *very specific* form, namely, it consists of subblocks representing matrices proportional to the *identity matrix* $`\mathrm{I}_{n_j}`$ of the dimension $`n_j`$ which repeat $`m_j`$ times along the rows and columns of the block $`\mathrm{D}_j`$. We can illustrate the structure of a certain block $`\mathrm{D}_j=\mathrm{D}`$ characterized by the numbers $`n_j=n`$, $`m_j=m`$ as follows
$$\mathrm{D}=\left(\begin{array}{cccc}\mu _{11}\mathrm{I}_n& \mu _{12}\mathrm{I}_n& \mathrm{}& \mu _{1m}\mathrm{I}_n\\ \mu _{21}\mathrm{I}_n& \mu _{22}\mathrm{I}_n& \mathrm{}& \mu _{2m}\mathrm{I}_n\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \mu _{m1}\mathrm{I}_n& \mu _{m2}\mathrm{I}_n& \mathrm{}& \mu _{mm}\mathrm{I}_n\end{array}\right),$$
(80)
where $`\mathrm{I}_n`$ is the $`n\times n`$ identity matrix.
In our case, the matrix $`\mathrm{H}`$ is the Jacobi matrix $`\mathrm{J}(t)`$ of the linearized system $`\ddot{𝜹}=\mathrm{J}(t)𝜹`$, $`G`$ is the symmetry group of a given bush B$`[G]`$, $`\mathrm{\Gamma }`$ is the mechanical representation of this group. Each block $`\mathrm{D}_j`$ generates an independent subsystem with $`n_jm_j`$ equations in the decomposition of the linearized system. However, each of these subsystems automatically splits into $`n_j`$ new subsystems consisting of $`m_j`$ differential equations, as a consequence of the specific structure of the block $`\mathrm{D}_j`$ (see Eq. (80)). Indeed, for example, if a certain $`\mathrm{D}`$-block for the matrix $`\mathrm{J}(t)`$ possesses the form ($`n_j=3`$, $`m_j=2`$):
$$\mathrm{J}_j(t)=\left(\begin{array}{cc}\mu _{11}\mathrm{I}_3& \mu _{12}\mathrm{I}_3\\ \mu _{21}\mathrm{I}_3& \mu _{22}\mathrm{I}_3\end{array}\right),$$
it is easy to check that we obtain the following three independent pairs of the equations from the system $`\ddot{𝜹}=\mathrm{J}(t)𝜹`$:
$$\begin{array}{c}\{\begin{array}{c}\ddot{\delta }_1=\mu _{11}\delta _1+\mu _{12}\delta _4,\hfill \\ \ddot{\delta }_4=\mu _{21}\delta _1+\mu _{22}\delta _4,\hfill \end{array}\hfill \\ \\ \{\begin{array}{c}\ddot{\delta }_2=\mu _{11}\delta _2+\mu _{12}\delta _5,\hfill \\ \ddot{\delta }_5=\mu _{21}\delta _2+\mu _{22}\delta _5,\hfill \end{array}\hfill \\ \\ \{\begin{array}{c}\ddot{\delta }_3=\mu _{11}\delta _3+\mu _{12}\delta _6,\hfill \\ \ddot{\delta }_6=\mu _{21}\delta _3+\mu _{22}\delta _6.\hfill \end{array}\hfill \end{array}$$
(81)
Note that the dimension of each subsystem (81) is equal to $`m_j=2`$, while the total number of these subsystems is equal to $`n_j=3`$.
#### IV.5.3 Irreducible representations and their basis vectors for the dihedral group
Hereafter, for simplicity, we will discuss only chains with an *even* number ($`N`$) of particles and illustrate the main ideas with the example $`N=12`$.
The symmetry of an $`N`$-particle (monoatomic) chain is completely described by the *dihedral* group $`G_0=D_N`$ which can be written as the union of two cosets with respect to its translational subgroup $`T_N=\{\widehat{e},\widehat{a},\widehat{a}^2,\mathrm{},\widehat{a}^{N1}\}`$:
$$D_N=T_NT_N\widehat{ı}.$$
(82)
Here $`\widehat{ı}`$ is the inversion relative to the center of the chain. The group $`D_N`$ is a non-Abelian group, since some of its elements do not commute with each other (for example, $`\widehat{ı}\widehat{a}=\widehat{a}^1\widehat{ı}`$). As a consequence, the number of classes of conjugate elements of this group is less than the total number ($`2N`$) of its elements and some irreps $`\mathrm{\Gamma }_j`$ are not one-dimensional. The irreps of the dihedral group $`D_N`$ can be obtained by the well-known induction procedure from those of its subgroup $`T_N`$. It turns out that for $`D_N`$ with *even* $`N`$ there are four one-dimensional irreps, while all the other ($`\frac{N}{2}1`$) irreps are two-dimensional. We discussed the construction of these irreps in FPU1 , where the following results were obtained.
Every irrep can be determined by two matrices $`\mathrm{M}_j(\widehat{a})`$, $`\mathrm{M}_j(\widehat{ı})`$ corresponding to its generators $`\widehat{a}`$ and $`\widehat{ı}`$, where $`j`$ is the number of this irrep. Four one-dimension irreps ($`j=1,2,3,4`$) are real and are determined by matrices <sup>8</sup><sup>8</sup>8All combinations of signs are allowed in (83).
$$\mathrm{M}_j(\widehat{a})=\pm 1,\mathrm{M}_j(\widehat{ı})=\pm 1.$$
(83)
All other irreps are two-dimensional and are determined by matrices
$$\mathrm{M}_j(\widehat{a})=\left(\begin{array}{cc}\mu _j& 0\\ 0& \overline{\mu }_j\end{array}\right),\mathrm{M}_j(\widehat{ı})=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),$$
with $`\mu _j=e^{\frac{2\pi ij}{N}}`$, $`\overline{\mu }_j=e^{\frac{2\pi ij}{N}}`$ ($`j0,N/2`$) <sup>9</sup><sup>9</sup>9For the values $`j=0`$ and $`j=N/2`$, two-dimensional representations turn out to be reducible and they decompose into two pairs of one-dimensional irreps listed in (83)..
Let us find the basis vectors of the irreducible representations of the dihedral group $`[\widehat{a}^m,\widehat{ı}]`$ for the case $`m=3`$ which corresponds to the bush B$`[\widehat{a}^m,\widehat{ı}]`$. Let $`\mathit{\varphi }`$ and $`𝝍`$ be the basis vectors of the two-dimensional invariant subspace corresponding to the irrep with the matrix $`\mathrm{M}(\widehat{g})=\left(\begin{array}{cc}\gamma & 0\\ 0& \overline{\gamma }\end{array}\right)`$, where $`\widehat{g}=\widehat{a}^m`$ is the translational generator of the dihedral group. They can be obtained from the equations $`\widehat{g}\mathit{\varphi }=\gamma \mathit{\varphi }`$ and $`\widehat{g}𝝍=\overline{\gamma }𝝍`$, respectively. For example, using Eq. (58) for the case $`m=3`$, we find $`\mathit{\varphi }=\{x,y,z|\gamma ^1x,\gamma ^1y,\gamma ^1z|\gamma ^2x,\gamma ^2y,\gamma ^2z|\mathrm{}\}`$, $`𝝍=\{\stackrel{~}{x},\stackrel{~}{y},\stackrel{~}{z}|\overline{\gamma }^1\stackrel{~}{x},\overline{\gamma }^1\stackrel{~}{y},\overline{\gamma }^1\stackrel{~}{z}|\overline{\gamma }^2\stackrel{~}{x},\overline{\gamma }^2\stackrel{~}{y},\overline{\gamma }^2\stackrel{~}{z}|\mathrm{}\}`$. Here ($`x,y,z`$) and ($`\stackrel{~}{x},\stackrel{~}{y},\stackrel{~}{z}`$) are arbitrary constants which these vectors depend on.
Taking into account the presence of the matrix $`\mathrm{M}(\widehat{ı})=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$ in every two-dimensional irrep, one can state that
$$\widehat{ı}\mathit{\varphi }=𝝍,\widehat{ı}𝝍=\mathit{\varphi }.$$
(84)
Because of these relations, there appear certain connections between the arbitrary constants ($`x,y,z`$) and ($`\stackrel{~}{x},\stackrel{~}{y},\stackrel{~}{z}`$). As a consequence, the basis vectors $`\mathit{\varphi }`$ and $`𝝍`$, for each two-dimensional irrep of the group $`[\widehat{a}^3,\widehat{ı}]`$, depends on only three arbitrary parameters: $`x`$, $`y`$ and $`z`$. In turn, this means that each two-dimensional irrep enters exactly *three* times into the decomposition of the mechanical representation of the dihedral group $`[\widehat{a}^3,\widehat{ı}]`$.
Unlike this, the one-dimensional irreps of the dihedral group $`[\widehat{a}^3,\widehat{ı}]`$ are contained in the mechanical representation less than $`3`$ times. Indeed, let us consider the basis vectors $`\mathit{\varphi }`$ and $`𝝍`$ of the one-dimensional irreps of the group $`[\widehat{a}^3]`$ determined by the matrices $`\mathrm{M}(a^3)=1`$ and $`\mathrm{M}(a^3)=1`$, respectively, for the case $`m=3`$, $`N=12`$:
$$\begin{array}{c}\mathit{\varphi }=\{x,y,z|x,y,z|x,y,z|x,y,z\},\hfill \\ 𝝍=\{x,y,z|x,y,z|x,y,z|x,y,z\}.\hfill \end{array}$$
These vectors can be obtained from Eq. (58) by letting $`\gamma =1`$ and $`\gamma =1`$. If the vector $`\mathit{\varphi }`$ is not only the basis vector of the irrep of the group $`[\widehat{a}^3]`$, but also is the basis vector of a certain one-dimensional irrep of the dihedral group $`[\widehat{a}^3,\widehat{ı}]`$, it must satisfy the equations $`\widehat{ı}\mathit{\varphi }=\mathit{\varphi }`$, for the irrep $`\mathrm{\Gamma }_1`$ ($`\mathrm{M}(\widehat{ı})=1`$) and $`\widehat{ı}\mathit{\varphi }=\mathit{\varphi }`$, for the irrep $`\mathrm{\Gamma }_2`$ ($`\mathrm{M}(\widehat{ı})=1`$). We obtain $`z=x`$, $`y=0`$ from the former equation, and $`z=x`$ for the latter equation. Thus
$$\begin{array}{c}\mathit{\varphi }[\mathrm{\Gamma }_1]=\{x,0,x|x,0,x|x,0,x|x,0,x\},\hfill \\ \mathit{\varphi }[\mathrm{\Gamma }_2]=\{x,y,x|x,y,x|x,y,x|x,y,x\}.\hfill \end{array}$$
In the same manner, we obtain the basis vectors $`𝝍[\mathrm{\Gamma }_3]`$ and $`𝝍[\mathrm{\Gamma }_4]`$ from the equations $`\widehat{ı}𝝍=𝝍`$ and $`\widehat{ı}𝝍=𝝍`$, respectively:
$$\begin{array}{c}𝝍[\mathrm{\Gamma }_3]=\{x,y,x|x,y,x|x,y,x|x,y,x\},\hfill \\ 𝝍[\mathrm{\Gamma }_4]=\{x,0,x|x,0,x|x,0,x|x,0,x\}.\hfill \end{array}$$
From the above results, we conclude that the irreps $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_4`$ are contained once, while the irreps $`\mathrm{\Gamma }_2`$ and $`\mathrm{\Gamma }_3`$ are contained twice in the decomposition of the mechanical representation for the considered chain.
The generalization of these results to the case of the dihedral group $`[\widehat{a}^m,\widehat{ı}]`$ with *arbitrary* $`m`$ is trivial.
The splitting scheme (78) for the bush B$`[\widehat{a}^3,\widehat{ı}]`$ for the monoatomic chain with $`N=12`$ particles can be now explained as follows. There are five irreps ($`\mathrm{\Gamma }_1`$, $`\mathrm{\Gamma }_2`$, $`\mathrm{\Gamma }_3`$, $`\mathrm{\Gamma }_4`$, $`\mathrm{\Gamma }_5`$) of the group $`[\widehat{a}^3,\widehat{ı}]\{\widehat{e},\widehat{a}^3,\widehat{a}^6,\widehat{a}^9|\widehat{ı},\widehat{ı}\widehat{a}^3,\widehat{ı}\widehat{a}^6,\widehat{ı}\widehat{a}^9\}D_4`$. As have just shown, the one-dimensional irreps $`\mathrm{\Gamma }_1`$ ($`n_1=1`$) and $`\mathrm{\Gamma }_4`$ ($`n_4=1`$) are contained once ($`m_1=1`$, $`m_4=1`$) in the decomposition of the mechanical representation $`\mathrm{\Gamma }`$ for our chain. On the other hand, the one-dimensional irreps $`\mathrm{\Gamma }_2`$ ($`n_2=1`$) and $`\mathrm{\Gamma }_3`$ ($`n_3=1`$) are contained twice ($`m_2=2`$, $`m_3=2`$) in $`\mathrm{\Gamma }`$, while two-dimensional irrep $`\mathrm{\Gamma }_5`$ ($`n_5=2`$) is contained thrice ($`m_5=3`$) in $`\mathrm{\Gamma }`$. The twelve variables $`\delta _j`$ from Eq. (78) are associated with the irreps of the group $`D_4`$ in the following manner:
$$\begin{array}{ccc}\delta _1\mathrm{\Gamma }_1(1);\hfill & & \\ \delta _2\mathrm{\Gamma }_2(1),\hfill & \delta _3\mathrm{\Gamma }_2(2);\hfill & \\ \delta _4\mathrm{\Gamma }_3(1),\hfill & \delta _5\mathrm{\Gamma }_3(2);\hfill & \\ \delta _6\mathrm{\Gamma }_4(1);\hfill & & \\ (\delta _7,\delta _8)\mathrm{\Gamma }_5(1),\hfill & (\delta _9,\delta _{10})\mathrm{\Gamma }_5(2),\hfill & (\delta _{11},\delta _{12})\mathrm{\Gamma }_5(3).\hfill \end{array}$$
(85)
Here, in parenthesis, we give the index of the copy of the irrep $`\mathrm{\Gamma }_i`$ (whose dimension is equal to $`n_i`$) in the decomposition of the mechanical representation $`\mathrm{\Gamma }`$. Note that the total number of such copies determines how many times $`m_i`$ the irrep $`\mathrm{\Gamma }_i`$ is contained in $`\mathrm{\Gamma }`$. On the other hand, as we already know, $`m_i`$ shows us the dimension of the subsystems $`L_j`$, while $`n_i`$ determines the total number of $`L_j`$ with the same dimension associated with $`\mathrm{\Gamma }_i`$. As a result, we obtain the splitting scheme (78).
The above discussed decomposition of the full linearized system $`\ddot{𝜹}=\mathrm{J}(t)𝜹`$ into independent subsystems $`L_j`$ of small dimensions permits one to analyze efficiently the stability of a given bush in the monoatomic chain with arbitrary large number of particles ($`N`$).
Using this idea, the stability diagrams for all the one-dimensional bushes in both FPU-$`\alpha `$ and FPU-$`\beta `$ chains were obtained in FPU2 . As an example, in Fig. 1, we reproduce the stability diagram for the bush B$`[\widehat{a}^3,\widehat{ı}]`$ for the FPU-$`\alpha `$ chain from that paper. In this diagram, each point $`(A,q)`$ determines a certain value of the bush mode amplitude $`A`$ and a certain value of the wave number $`q=\frac{2\pi j}{N}`$ that is associated with the index $`j`$ of a fixed mode. The black points $`(A,q)`$ correspond to the case where the mode $`j=q\frac{N}{2\pi }`$ becomes excited because of its parametric interaction with the mode of the bush B$`[\widehat{a}^3,\widehat{ı}]`$. The white color denotes the opposite case: the corresponding mode $`j`$, being zero at the initial instant, continues to be zero in spite of its interaction with the considered bush. Such a diagram allows one to study stability of one-dimensional bushes not only for finite $`N`$, but also for the case $`N\mathrm{}`$ (some more details can be found in FPU2 ).
## V Some additional examples
### V.1 Stability of the $`\pi `$-mode in the FPU-$`\alpha `$ chain
The symmetry group of the $`\pi `$-mode $`𝑿(t)=\{A(t),A(t),A(t),A(t),A(t),A(t),\mathrm{}\}`$ in the FPU-$`\alpha `$ chain is $`[\widehat{a}^2,\widehat{ı}]`$. Using results of the previous sections, one can deduce for $`N`$ divisible by $`4`$ that the linearized system in the vicinity of this mode splits into four individual equations and a number of two-dimensional systems of differential equations <sup>10</sup><sup>10</sup>10If $`N`$ is an even number, but $`Nmod40`$, one obtains two equations of harmonic oscillators and a number of pairs of coupled equations.. The first pair of individual equations represents two independent harmonic oscillators, while the second pair represents two Mathieu equations. All the other systems represent pairs of coupled equations:
$$\begin{array}{c}\ddot{\delta }_j+4\mathrm{sin}^2\left(\frac{\pi j}{N}\right)\delta _j=\eta \mathrm{sin}\left(\frac{2\pi j}{N}\right)\delta _{\frac{N}{2}j}\mathrm{cos}(2t),\hfill \\ \ddot{\delta }_{\frac{N}{2}j}+4\mathrm{cos}^2\left(\frac{\pi j}{N}\right)\delta _{\frac{N}{2}j}=\eta \mathrm{sin}\left(\frac{2\pi j}{N}\right)\delta _j\mathrm{cos}(2t),\hfill \end{array}$$
(86)
where $`\eta =\frac{8A}{\sqrt{N}}`$, $`j=1,2,\mathrm{},\frac{N}{4}1`$. In (86), $`A`$ is the amplitude of the $`\pi `$-mode, i.e. $`A=\mathrm{max}|A(t)|`$. Eqs. (86) can be rewritten as follows:
$$\begin{array}{c}\ddot{x}+4\mathrm{sin}^2\left(\frac{q}{2}\right)x=\eta \mathrm{sin}(q)y\mathrm{cos}(2t),\hfill \\ \ddot{y}+4\mathrm{cos}^2\left(\frac{q}{2}\right)y=\eta \mathrm{sin}(q)x\mathrm{cos}(2t),\hfill \end{array}$$
(87)
where $`q=\frac{2\pi j}{N}`$ is the wave number, $`x(t)=\delta _j(t)`$, $`y(t)=\delta _{\frac{N}{2}j}(t)`$. These results were obtained and discussed in FPU1 using a different (in comparison with the present paper) method. Indeed, there we obtained *exact* equations for the FPU-$`\alpha `$ chain in the *modal* space and only then linearized them near the $`\pi `$-mode.
As it was discussed in FPU1 , the system (87) turns out to be rather remarkable: the $`\pi `$-mode (or one-dimensional bush B$`[\widehat{a}^2,\widehat{ı}]`$) loses its stability *simultaneously* with respect to interaction with all the other modes. In other words, the threshold $`\eta _c`$ for the loss of stability for the bush B$`[\widehat{a}^2,\widehat{ı}]`$ turns out to be *the same* for all the values of $`q`$, i.e. for all the subsystems (86):
$$\eta _c=2.42332\mathrm{}$$
(88)
This property of the linearized system (87) was discussed in more detail in FPU1 . The stability diagram for the bush B$`[\widehat{a}^2,\widehat{ı}]`$ can be found in FPU2 .
### V.2 Stability of the $`\pi `$-mode in the FPU-$`\beta `$ chain
This case is of a particular interest since the linearized system for the $`\pi `$-mode in the FPU-$`\beta `$ chain possesses a *higher symmetry* compared to that in the FPU-$`\alpha `$ chain. Indeed, for the FPU-$`\beta `$ model, the interparticle potential is an *even* function, and some additional symmetry elements of the dynamical equations appear as a consequence of this fact.
Let us introduce an operator $`\widehat{u}`$ that changes signs of the displacements of all the particles:
$$\widehat{u}\{x_1,x_2,\mathrm{},x_N\}=\{x_1,x_2,\mathrm{},x_N\}.$$
(89)
This operator leaves the FPU-$`\beta `$ Hamiltonian unchanged because of its evenness.
Therefore, the symmetry group of the FPU-$`\beta `$ dynamical equations turns out to be a *supergroup* $`G_0^{}`$ with respect to the symmetry group $`G_0=[\widehat{a},\widehat{ı}]`$ of the FPU-$`\alpha `$ chain:
$$G_0^{}=G_0G_0\widehat{u}$$
(90)
(note that $`\widehat{u}^2=\widehat{e}`$).
As a result, one can classify bushes of modes in the FPU-$`\beta `$ chain by subgroups of the group $`G_0^{}`$ rather than by subgroups of the $`G_0`$. The invariant manifolds (bushes of modes) in the FPU-$`\beta `$ model with respect to the group $`G_0^{}`$ were found by Rink in BR . Some additional details of this problem were discussed in our paper FPU2 (in particular, the dynamical equations and stability of these bushes of modes).
Considering $`G_0^{}`$ from (90) as the parent group, we discover that the $`\pi `$-mode in the FPU-$`\beta `$ chain must be characterized by the group $`G=[\widehat{a}^2,\widehat{ı},\widehat{a}\widehat{u}]`$, unlike the group $`[\widehat{a}^2,\widehat{ı}]`$ characterizing the $`\pi `$-mode in the FPU-$`\alpha `$ chain with the parent group $`[\widehat{a},\widehat{ı}]`$.
The third generator <sup>11</sup><sup>11</sup>11Obviously, the operators $`\widehat{a}`$ and $`\widehat{u}`$ commute with each other. $`\widehat{a}\widehat{u}\widehat{u}\widehat{a}`$ of the group $`G=[\widehat{a}^2,\widehat{ı},\widehat{a}\widehat{u}]`$ acts on the configuration vector $`𝑿`$ as follows:
$$\widehat{a}\widehat{u}𝑿=\{x_N,x_1,x_2,\mathrm{},x_{N1}\}.$$
(91)
Above, we specified the group $`G`$ by *three* generators, but it can be determined by only *two* generators. Indeed, if
$$\widehat{p}=\widehat{a}\widehat{u},$$
(92)
then $`\widehat{p}^2=\widehat{a}\widehat{u}\widehat{a}\widehat{u}=\widehat{a}^2\widehat{u}^2=\widehat{a}^2`$. Thus, the first generator in the list $`G=[\widehat{a}^2,\widehat{ı},\widehat{a}\widehat{u}]`$ is simply the square of the third generator $`\widehat{p}=\widehat{a}\widehat{u}`$.
For simplicity, let us consider the case $`N=6`$. Then $`[\widehat{a}^2]=\{\widehat{e},\widehat{a}^2,\widehat{a}^4\}`$, ($`\widehat{a}^6=\widehat{e}`$) is a cyclic group of the order $`3`$ and the full group $`G`$ contains $`322=12`$ elements. With the aid of the operators $`\widehat{p}=\widehat{a}\widehat{u}`$ and $`\widehat{q}=\widehat{ı}`$, we can obtain all the elements of the group
$$G=[\widehat{a}^2,\widehat{ı},\widehat{a}\widehat{u}][\widehat{p},\widehat{q}]$$
(93)
as follows:
$$G=T_6T_6\widehat{q},$$
(94)
where $`T_6`$ is the cyclic group of the order $`6`$:
$$T_6=\{\widehat{e},\widehat{p},\widehat{p}^2,\widehat{p}^3,\widehat{p}^4,\widehat{p}^5\}.$$
(95)
The following generating relations, fully determining the group $`G=[\widehat{p},\widehat{q}]`$, can be obtained:
$$\widehat{p}^6=\widehat{e},\widehat{q}^2=\widehat{e},\widehat{q}\widehat{p}=\widehat{p}^5\widehat{q}.$$
(96)
From these relations, one can see that $`G`$ is the *dihedral* group $`D_6`$ (by the way, it is isomorphic to the point groups $`C_{6v}`$, $`D_{3h}`$ and $`D_{3d}`$, as well).
The irreducible representations of the dihedral group were discussed in FPU1 (see also Sec. IV.5 of the present paper). There are four one-dimensional and $`\left(\frac{N}{2}1\right)`$ two-dimensional irreps of the group $`D_N`$ with an even index $`N`$. For simplicity, let us discuss the case $`N=6`$ only (the generalization to arbitrary values of $`N`$ turns out to be trivial).
All irreps of the group $`D_6`$ can be constructed from the irreps of its subgroup $`T_6`$ (see (94)) with the aid of the induction procedure. As a result, we obtain the following irreps of the group $`G=[\widehat{p},\widehat{q}]`$ presented in Table 3. In this table, each irrep $`\mathrm{\Gamma }_j`$ ($`j=1,2,\mathrm{},6`$) is determined by the matrices $`\mathrm{M}(\widehat{p})`$ and $`\mathrm{M}(\widehat{q})`$, corresponding to the generators $`\widehat{p}`$ and $`\widehat{q}`$ of the group $`D_6`$.
All the invariant subspaces of the configuration space of the chain with $`N=6`$ particles, corresponding to these irreps, and their basis vectors can be obtained as follows. Let us find the basis vector $`\mathit{\varphi }`$ of a certain one-dimensional irrep $`\mathrm{\Gamma }`$ determined by the matrices $`\mathrm{M}(\widehat{p})=\gamma `$ and $`\mathrm{M}(\widehat{q})=\delta `$ from Table 3 ($`\gamma =\pm 1`$, $`\delta =\pm 1`$). Thus, the vector $`\mathit{\varphi }`$ must satisfy the relations:
$`\widehat{p}\mathit{\varphi }=\gamma \mathit{\varphi },`$ (97a)
$`\widehat{q}\mathit{\varphi }=\delta \mathit{\varphi }.`$ (97b)
From (97a), we obtain
$$\mathit{\varphi }=\{x,\gamma ^5x,\gamma ^4x,\gamma ^3x,\gamma ^2x,\gamma x\},$$
(98)
where $`x`$ is an arbitrary constant. For $`\gamma =1`$, the basis vector $`\mathit{\varphi }`$ from (98) transforms to
$$\mathit{\varphi }_1=\{x,x,x,x,x,x\},$$
(99)
while for $`\gamma =1`$ it transforms to
$$\mathit{\varphi }_2=\{x,x,x,x,x,x\}.$$
(100)
Now we must also demand (97b) to hold. Remembering that $`\widehat{q}=\widehat{ı}`$, and, therefore, that $`\widehat{q}\mathit{\varphi }=\{x_6,x_5,x_4,x_3,x_2,x_1\}`$, we obtain
$$\begin{array}{c}\widehat{q}\mathit{\varphi }_1=\{x,x,x,x,x,x\}(1)\mathit{\varphi }_1,\hfill \\ \widehat{q}\mathit{\varphi }_2=\{x,x,x,x,x,x\}(1)\mathit{\varphi }_2.\hfill \end{array}$$
This means that $`\mathit{\varphi }_1`$ from (99) is the basis vector of the irrep $`\mathrm{\Gamma }_1`$ ($`\gamma =1`$, $`\delta =1`$, see Eqs. (97)), while $`\mathit{\varphi }_2`$ from (100) is the basis vector of the irrep $`\mathrm{\Gamma }_4`$ ($`\gamma =1`$, $`\delta =1`$). Thus, each irrep $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_4`$ is contained once in the decomposition of the mechanical representations.
On the other hand, the basis vectors of irreps $`\mathrm{\Gamma }_2`$ ($`\gamma =1`$, $`\delta =1`$) and $`\mathrm{\Gamma }_3`$ ($`\gamma =1`$, $`\delta =1`$) are equal to zero. Indeed, demanding $`\widehat{q}\mathit{\varphi }_1=(1)\mathit{\varphi }_1`$ ($`\delta =1`$) one concludes that $`x=x`$ and, therefore, $`x=0`$ and $`\mathit{\varphi }_10`$. The same result originates also from the equation $`\widehat{q}\mathit{\varphi }_2=(1)\mathit{\varphi }_2`$ ($`\delta =1`$), namely, $`\mathit{\varphi }_20`$. Thus, we must conclude that the irreps $`\mathrm{\Gamma }_2`$ and $`\mathrm{\Gamma }_3`$ *are not contained* in the decomposition of the mechanical representation of our chain into irreducible representations (the corresponding invariant subspaces turn out to be null spaces).
There are two basis vectors, $`\mathit{\varphi }`$ and $`𝝍`$, for each two-dimensional irreps from Table 3. For the irrep $`\mathrm{\Gamma }_5`$ the following relations must hold
$`\widehat{p}\mathit{\varphi }=\mu \mathit{\varphi },`$ (101a)
$`\widehat{p}𝝍=\overline{\mu }𝝍,`$ (101b)
$`\widehat{q}\mathit{\varphi }=𝝍,`$ (101c)
$`\widehat{q}𝝍=\mathit{\varphi }.`$ (101d)
Comparing (101a) and (101b) with (97a), we can write
$$\mathit{\varphi }=\{x,\mu ^5x,\mu ^4x,\mu ^3x,\mu ^2x,\mu x\},$$
(102)
$$𝝍=\{y,\overline{\mu }^5y,\overline{\mu }^4y,\overline{\mu }^3y,\overline{\mu }^2y,\overline{\mu }y\},$$
(103)
where $`x`$ and $`y`$ are two *different* constants.
Then, from (101c) or (101d), we obtain a certain relations between the constants $`x`$ and $`y`$:
$$y=\mu x.$$
(104)
Substituting this value of $`y`$ into (103) we, finally, conclude that the irrep $`\mathrm{\Gamma }_5`$ is contained only once in the decomposition of the mechanical representation of our FPU-$`\beta `$ chain. Indeed, the basis vectors $`\mathit{\varphi }`$, $`𝝍`$ turn out to depend on only one constant $`x`$ which can be determined from the normalization condition. The same conclusion is valid also for the second irrep $`\mathrm{\Gamma }_6`$ of the group $`D_6`$.
The generalization of this conclusion to the case of the FPU-$`\beta `$ chain with an arbitrary even number $`N=2n`$ of particles can be achieved trivially. Namely, each two-dimensional irrep of the group $`D_{2n}`$ is contained only *once* in the mechanical representation of this group for the FPU-$`\beta `$ chain. One-dimensional irreps are contained in the mechanical representation once or not at all. Remembering the above discussion about the application of the Wigner theorem to splitting the linearized system $`\ddot{𝜹}=\mathrm{J}(t)𝜹`$, we discover that for the FPU-$`\beta `$ chain (unlike FPU-$`\alpha `$ chain) this system is decomposed into individual differential equations of the second order, i.e. we have the *complete splitting* in this case.
The following remark should be done to avoid a possible misunderstanding. The form of the $`\pi `$-mode
$$𝑿(t)=\{A(t),A(t),A(t),A(t),A(t),A(t)\},$$
(105)
is one and the same <sup>12</sup><sup>12</sup>12Note that $`A(t)`$ in Eq. (105) are *different* functions of time for the FPU-$`\alpha `$ and FPU-$`\beta `$ chains (see below). for both FPU-$`\alpha `$ and FPU-$`\beta `$ chains. Therefore, its symmetry group can be written as $`[\widehat{a}^2,\widehat{ı},\widehat{a}\widehat{u}]`$ rather than $`[\widehat{a}^2,\widehat{ı}]`$ since the operator $`\widehat{a}\widehat{u}`$ does not change the pattern (105) not only for the FPU-$`\beta `$ chain, but also for the FPU-$`\alpha `$ chain.
Nevertheless, the linearized system $`\ddot{𝜹}=\mathrm{J}(t)𝜹`$, for the FPU-$`\beta `$ chain, is invariant with respect to the operator $`\widehat{a}\widehat{u}`$, while that for the FPU-$`\alpha `$ chain is not invariant under the action of $`\widehat{a}\widehat{u}`$. Indeed, the linearized (near the $`\pi `$-mode) system for the FPU-$`\alpha `$ chain reads:
$$\begin{array}{c}\ddot{\delta }_1=[\delta _22\delta _1+\delta _6]4A(t)[\delta _2\delta _6],\hfill \\ \ddot{\delta }_2=[\delta _32\delta _2+\delta _1]+4A(t)[\delta _3\delta _1],\hfill \\ \ddot{\delta }_3=[\delta _42\delta _3+\delta _2]4A(t)[\delta _4\delta _2],\hfill \\ \ddot{\delta }_4=[\delta _52\delta _4+\delta _3]+4A(t)[\delta _5\delta _3],\hfill \\ \ddot{\delta }_5=[\delta _62\delta _5+\delta _4]4A(t)[\delta _6\delta _4],\hfill \\ \ddot{\delta }_6=[\delta _12\delta _6+\delta _5]+4A(t)[\delta _1\delta _5].\hfill \end{array}$$
(106)
On the other hand, for the FPU-$`\beta `$ chain, the linearized system reads:
$$\begin{array}{c}\ddot{\delta }_1=[\delta _22\delta _1+\delta _6][1+12A^2(t)],\hfill \\ \ddot{\delta }_2=[\delta _32\delta _2+\delta _1][1+12A^2(t)],\hfill \\ \ddot{\delta }_3=[\delta _42\delta _3+\delta _2][1+12A^2(t)],\hfill \\ \ddot{\delta }_4=[\delta _52\delta _4+\delta _3][1+12A^2(t)],\hfill \\ \ddot{\delta }_5=[\delta _62\delta _5+\delta _4][1+12A^2(t)],\hfill \\ \ddot{\delta }_6=[\delta _12\delta _6+\delta _5][1+12A^2(t)].\hfill \end{array}$$
(107)
The operator $`\widehat{a}\widehat{u}`$ acts on the vector $`𝜹`$ as a follows:
$$\widehat{a}\widehat{u}𝜹\widehat{a}\widehat{u}\{\delta _1,\delta _2,\delta _3,\delta _4,\delta _5,\delta _6\}=\{\delta _6,\delta _1,\delta _2,\delta _3,\delta _4,\delta _5\}.$$
(108)
The substitution
$$\begin{array}{c}\delta _1\delta _6,\delta _2\delta _1,\delta _3\delta _2,\hfill \\ \delta _4\delta _3,\delta _5\delta _4,\delta _6\delta _5,\hfill \end{array}$$
(109)
transforms, for example, the first equation of the system (106) to the form
$$\ddot{\delta }_6=[\delta _12\delta _6+\delta _5]4A(t)[\delta _1\delta _5],$$
which differs by the sign in front of $`4A(t)`$ from the sixth equation of the system (106). Contrariwise, the transformation (109) leaves the system (107) for the FPU-$`\beta `$ chain unchanged (it leads only to some transpositions of the individual equation in (107)).
What does it mean? Each bush (the $`\pi `$-mode, in our case) is associated with a certain subgroup of the parent group $`G_0`$, the symmetry group of the original nonlinear *dynamical equations* of the considered mechanical system. The operator $`\widehat{u}`$ does not belong to the symmetry group $`G_0`$ of the FPU-$`\alpha `$ chain (unlike the case of the FPU-$`\beta `$ chain!) and, therefore, there are no subgroups of the group $`G_0`$ whose elements contain $`\widehat{u}`$. Precisely this fact can explain why we must not take into consideration the operator $`\widehat{a}\widehat{u}`$ for the case of the FPU-$`\alpha `$ chain, even through this operator does not change the pattern (105) of the $`\pi `$-mode.
The function $`A(t)`$ from the displacement pattern (105) of the $`\pi `$-mode is determined by the dynamical equation of the one-dimensional bush B$`[\widehat{a}^2,\widehat{ı}]`$. As it was already discussed, this equation can be obtained by substituting the vector $`𝑿(t)`$ from Eq. (105) into the nonlinear dynamical equations of the FPU chain. For the FPU-$`\alpha `$ model this equation turns out to be the equation of a harmonic oscillator $`\ddot{A}+4A=0`$, while for the FPU-$`\beta `$ model it is the Duffing equation $`\ddot{A}+4A+16A^3=0`$.
Thus, in both cases, the exact expression for the function $`A(t)`$ can be found. For studying the stability of the $`\pi `$-mode, we must substitute the corresponding expression for $`A(t)`$ into the linearized system $`\ddot{𝜹}=\mathrm{J}(t)𝜹`$. As we saw in the present section, the linearized system splits into individual equations for the FPU-$`\beta `$ chain, while for the FPU-$`\alpha `$ chain it can be decomposed into four individual equations and $`(N4)/2`$ pairs of differential equations. All these equations turn out to be equations of the second order with time-periodic coefficients depending on the function $`A(t)`$ FPU1 ; FPU2 .
Let us note that stability of the $`\pi `$-mode in the FPU-$`\alpha `$ and FPU-$`\beta `$ chains was investigated by different methods in a large number of papers (see, for example, \[Bud, ; Sand, ; Flach, , FPU1, ; FPU2, ; PR, ; Yosh, ; Shin, ; CLL, ; AntiFPU, \]), but, to our best understanding, the influence of symmetry of these mechanical models on stability analysis was not discussed. Unlike the above cited works, the bush stability analysis presented in this paper based on the symmetry-related arguments only. Therefore, our conclusion about the difference in splitting scheme of linearized systems for the $`\pi `$-mode in the FPU-$`\alpha `$ and FPU-$`\beta `$ chains can be automatically extended to all the other nonlinear chains with the same symmetry characteristics. In particular, we can conclude that the splitting of the linearized system into individual equations can be performed not only for the FPU-$`\beta `$ chain, but for every chain with an *even* potential of the interparticle interaction. In contrast, it is impossible for the chains with *arbitrary* potential, and the FPU-$`\alpha `$ model is a simple illustration of this proposition.
We would like to focus on the paper Yosh , where some analytical results were obtained for the stability of the $`\pi `$-mode in the nonlinear chains with a *general even* form of the interparticle interaction potential. The author of Yosh succeeded in his analysis thanks to the decomposition of the linearized system into individual equations, but such analysis cannot be extended to the FPU-$`\alpha `$ chain precisely because the potential of this model is *odd*.
### V.3 Stability of two-dimensional bush B$`[\widehat{a}^4,\widehat{ı}]`$ in the FPU-$`\alpha `$ chain
Now, let us consider the stability of the two-dimensional bush B$`[\widehat{a}^4,\widehat{ı}]`$ that can be determined by the displacement pattern FPU2
$$𝑿(t)=\{A(t),B(t),B(t),A(t)|A(t),B(t),B(t),A(t)|A(t),B(t),B(t),A(t)\}$$
(110)
(for simplicity, we start with the case $`N=12`$).
The symmetry group of this bush is the dihedral group $`D_3`$ with translational subgroup
$$T_3=\{\widehat{e},\widehat{a}^4,\widehat{a}^8\}(\widehat{a}^{12}=\widehat{e}).$$
(111)
This non-Abelian group ($`\widehat{ı}\widehat{a}^4=\widehat{a}^8\widehat{ı}`$) consists of six elements determined by the equation
$$D_3=T_3T_3\widehat{ı},$$
(112)
and possesses the following three irreps presented in Table 4 (we give there the two-dimensional irrep not only in the complex form, but in the real form also).
The basis vectors of the one-dimensional irreps of the group $`T_3`$ (111) can be found from the relation
$$\widehat{p}\mathit{\varphi }=\gamma \mathit{\varphi },$$
(113)
where $`\gamma =1,\mu ,\overline{\mu }`$. From this equation we obtain
$`\widehat{p}\mathit{\varphi }=\{x_9,x_{10},x_{11},x_{12}|x_1,x_2,x_3,x_4|x_5,x_6,x_7,x_8\}=`$
$`\gamma \{x_1,x_2,x_3,x_4|x_5,x_6,x_7,x_8|x_9,x_{10},x_{11},x_{12}\}`$ (114)
and then
$$\mathit{\varphi }=\{x_1,x_2,x_3,x_4|\overline{\gamma }x_1,\overline{\gamma }x_2,\overline{\gamma }x_3,\overline{\gamma }x_4|\gamma x_1,\gamma x_2,\gamma x_3,\gamma x_4\}.$$
(115)
This basis vector depends on four arbitrary constants ($`x_1`$, $`x_2`$, $`x_3`$, $`x_4`$) and, therefore, it determines a four-dimensional subspace invariant under the translational group $`T_3`$ (111), associated with the irrep $`\mathrm{\Gamma }`$ defined by the one-dimensional matrix $`\mathrm{M}(\widehat{p})=\gamma `$.
The basis vectors of one-dimensional irreps of the whole group $`D_3`$ from (112) can be obtained with the aid of the equation $`\widehat{ı}\mathit{\varphi }=\mathit{\varphi }`$ (for the irrep $`\mathrm{\Gamma }_1`$) and with the aid of the equation $`\widehat{ı}\mathit{\varphi }=\mathit{\varphi }`$ (for the irrep $`\mathrm{\Gamma }_2`$), where $`\mathit{\varphi }=\{x_1,x_2,x_3,x_4|x_1,x_2,x_3,x_4|x_1,x_2,x_3,x_4\}`$ since $`\gamma =1`$ for these both irreps (see Eq. (115)) <sup>13</sup><sup>13</sup>13Here $`\gamma `$ is the one-dimensional matrix associated with the generator $`\widehat{p}`$ (see Table 4).. From the equation $`\widehat{ı}\mathit{\varphi }=\mathit{\varphi }`$, we obtain
$$\mathit{\varphi }[\mathrm{\Gamma }_1]=\{x_1,x_2,x_2,x_1|x_1,x_2,x_2,x_1|x_1,x_2,x_2,x_1\}.$$
(116)
This result means that the irrep $`\mathrm{\Gamma }_1`$ is contained twice in the decomposition of the mechanical representation of the considered chain. Analogously, for the case of the irrep $`\mathrm{\Gamma }_2`$, we obtain from the equation $`\widehat{ı}\mathit{\varphi }=\mathit{\varphi }`$:
$$\mathit{\varphi }[\mathrm{\Gamma }_2]=\{x_1,x_2,x_2,x_1|x_1,x_2,x_2,x_1|x_1,x_2,x_2,x_1\}.$$
(117)
This vector also determines the two-dimensional subspace and, therefore, the irrep $`\mathrm{\Gamma }_2`$ of the group $`D_3`$ is contained twice in the decomposition of the mechanical representation of our chain.
Now let us consider the basis vectors of the two-dimensional irrep $`\mathrm{\Gamma }_3`$ from Table 4. The following relations must hold (for the complex form of this irrep)
$`\widehat{p}\mathit{\varphi }=\mu \mathit{\varphi },`$ (118a)
$`\widehat{p}𝝍=\overline{\mu }𝝍,`$ (118b)
$`\widehat{ı}\mathit{\varphi }=𝝍,`$ (118c)
$`\widehat{ı}𝝍=\mathit{\varphi },`$ (118d)
where $`\mathit{\varphi }`$ and $`𝝍`$ are two basis vectors of the irrep $`\mathrm{\Gamma }_3`$. Comparing (118a) and (118b) with (115) and letting $`\gamma =\mu `$ and $`\gamma =\overline{\mu }`$, respectively, we obtain:
$$\mathit{\varphi }=\{x_1,x_2,x_3,x_4|\overline{\mu }x_1,\overline{\mu }x_2,\overline{\mu }x_3,\overline{\mu }x_4|\overline{\mu }^2x_1,\overline{\mu }^2x_2,\overline{\mu }^2x_3,\overline{\mu }^2x_4\}$$
(119)
and
$$𝝍=\{y_1,y_2,y_3,y_4|\mu y_1,\mu y_2,\mu y_3,\mu y_4|\mu ^2y_1,\mu ^2y_2,\mu ^2y_3,\mu ^2y_4\}.$$
(120)
On the other hand, the both equations (118c) and (118d) lead to the same relation $`\widehat{ı}\mathit{\varphi }=𝝍`$. Then, from Eqs. (119), (120), we obtain $`y_1=\mu x_4`$, $`y_2=\mu x_3`$, $`y_3=\mu x_2`$, $`y_4=\mu x_1`$. Therefore, the final forms of the basis vectors of the irrep $`\mathrm{\Gamma }_3`$ are:
$$\mathit{\varphi }[\mathrm{\Gamma }_3]=\{x_1,x_2,x_3,x_4|\overline{\mu }x_1,\overline{\mu }x_2,\overline{\mu }x_3,\overline{\mu }x_4|\overline{\mu }^2x_1,\overline{\mu }^2x_2,\overline{\mu }^2x_3,\overline{\mu }^2x_4\},$$
(121)
$$𝝍[\mathrm{\Gamma }_3]=\{\mu x_4,\mu x_3,\mu x_2,\mu x_1|\mu ^2x_4,\mu ^2x_3,\mu ^2x_2,\mu ^2x_1,|x_4,x_3,x_2,x_1\}.$$
(122)
Each of these vectors depends on $`4`$ arbitrary parameters ($`x_1`$, $`x_2`$, $`x_3`$, $`x_4`$) and, therefore, one can construct four independent pairs of the basis vectors — $`\mathit{\varphi }_j[\mathrm{\Gamma }_3]`$, $`𝝍_j[\mathrm{\Gamma }_3]`$ ($`j=1,2,3,4`$) — of the irrep $`\mathrm{\Gamma }_3`$ <sup>14</sup><sup>14</sup>14It can be done by letting $`[x_1=1,x_2=0,x_3=0,x_4=0]`$, $`[x_1=0,x_2=1,x_3=0,x_4=0]`$, etc.. In turn, this means that the irrep $`\mathrm{\Gamma }_3`$ is contained four times in the decomposition of the mechanical representation of the chain with $`N=12`$ particles.
Taking into account the above results (see (116), (117), (121), (122)), we can conclude that in the case $`N=12`$, the linearized system $`\ddot{𝜹}=\mathrm{J}(t)𝜹`$ for studying the stability of the bush B$`[\widehat{a}^4,\widehat{ı}]`$ splits into two $`2\times 2`$ and two $`4\times 4`$ independent systems of differential equations of the second order. In the real form, these systems can be written as follows:
$$\begin{array}{c}\{\begin{array}{c}\ddot{\nu }_1+K_3(t)\nu _1=K_1(t)\nu _2,\hfill \\ \ddot{\nu }_2+K_4(t)\nu _2=K_1(t)\nu _1,\hfill \end{array}\hfill \\ \\ \{\begin{array}{c}\ddot{\nu }_3+K_1(t)\nu _3=K_1(t)\nu _4,\hfill \\ \ddot{\nu }_4+K_1(t)\nu _4=K_1(t)\nu _3.\hfill \end{array}\hfill \\ \\ \{\begin{array}{c}\ddot{\nu }_5+K_6(t)\nu _5=K_1(t)\nu _6+K_2(t)\nu _7,\hfill \\ \ddot{\nu }_6+K_3(t)\nu _6=K_1(t)\nu _5,\hfill \\ \ddot{\nu }_7+K_5(t)\nu _7=K_1(t)\nu _8+K_2(t)\nu _5,\hfill \\ \ddot{\nu }_8+K_1(t)\nu _8=K_1(t)\nu _7,\hfill \end{array}\hfill \\ \\ \{\begin{array}{c}\ddot{\nu }_9+K_6(t)\nu _9=K_1(t)\nu _{10}+K_2(t)\nu _{11},\hfill \\ \ddot{\nu }_{10}+K_3(t)\nu _{10}=K_1(t)\nu _9,\hfill \\ \ddot{\nu }_{11}+K_5(t)\nu _{11}=K_1(t)\nu _{12}+K_2(t)\nu _9,\hfill \\ \ddot{\nu }_{12}+K_1(t)\nu _{12}=K_1(t)\nu _{11},\hfill \end{array}\hfill \end{array}$$
(123)
where
$$\begin{array}{c}K_1(t)=12A(t)+2B(t),\hfill \\ K_2(t)=\sqrt{3}[\frac{1}{2}+2A(t)],\hfill \\ K_3(t)=3+6A(t)+2B(t),\hfill \\ K_4(t)=32A(t)6B(t),\hfill \\ K_5(t)=\frac{5}{2}2A(t)4B(t),\hfill \\ K_6(t)=\frac{3}{2}2A(t).\hfill \end{array}$$
The generalization to the case of an arbitrary $`N`$ (note that $`Nmod4=0`$ must hold!) is trivial: each one-dimensional irrep is contained twice in the decomposition of the mechanical representation of the $`N`$-particle chain, while each two-dimensional irrep <sup>15</sup><sup>15</sup>15The number of two-dimensional irreps of the bush symmetry group $`[\widehat{a}^4,\widehat{ı}]=D_{\frac{N}{4}}`$ increases with increasing $`N`$. is contained four times in it.
In conclusion, in Fig. 2 we reproduce the stability diagram for the two-dimensional bush B$`[\widehat{a}^4,\widehat{ı}]`$ from our paper FPU2 . This diagram corresponds to the FPU-$`\alpha `$ chain with $`N=12`$ particles. It represents a planar section of the four-dimensional stability domain in the space of the initial conditions $`\nu _1(0)`$, $`\nu _2(0)`$, $`\dot{\nu }_1(0)`$, $`\dot{\nu }_2(0)`$, where $`\nu _1(t)`$ and $`\nu _2(t)`$ are two modes of the considered bush.
In producing Fig. 2, we specify $`\dot{\nu }_1(0)=0`$, $`\dot{\nu }_2(0)=0`$, and change $`\nu _1(0)`$, $`\nu _2(0)`$ in some interval near their zero values. The stability domain, resembling a beetle, is drown in *black* color in the plain $`\nu _1(0)`$, $`\nu _2(0)`$. The bush B$`[\widehat{a}^4,\widehat{ı}]`$ losses its stability (and transforms into another bush of higher dimension), when we cross the boundary of the black region in any direction. From Fig. 2 it is obvious, how nontrivial stability domain for a bush of modes can be.
A detailed description of the stability domains for one-dimensional and two-dimensional bushes of modes in both FPU-$`\alpha `$ and FPU-$`\beta `$ chains can be found in FPU2 .
## VI Conclusion
All the exact dynamical regimes in $`N`$-particle mechanical system with discrete symmetry can be classified by the *subgroups* $`G_j`$ of the parent group $`G_0`$, i.e. the symmetry group of its equations of motion. Actually, each subgroup $`G_j`$ singles out a certain *invariant manifold* which, being decomposed into the basis vectors of the irreducible representations of the group $`G_0`$, is termed as a “bush of modes” DAN1 ; DAN2 ; PhysD .
The bush B$`[G_j]`$, representing an $`n`$-dimensional vibrational regime, can be considered as a dynamical object characterized by its displacement pattern of all the particles from their equilibrium positions, by the appropriate dynamical equations and the domain of the stability. One-dimensional bushes are symmetry-determined similar nonlinear normal modes introduced by Rosenberg Ros (see also IntJ ). For Hamiltonian systems, the energy of the initial excitation turns out to be “trapped” in the bush and this is a phenomenon of energy localization in the modal space.
The different aspects of the bush theory were developed in DAN1 ; DAN2 ; PhysD ; IntJ ; ENOC ; C60 ; Octa . Bushes of vibrational modes (invariant manifolds) in the FPU chains were discussed in FPU1 ; FPU2 ; BR ; PR ; Shin .
The stability analysis of a given bush B$`[G]`$ reduces to studying the linearized (in the vicinity of the bush) dynamical equations $`\ddot{𝜹}=\mathrm{J}(t)𝜹`$. In the present paper, we prove (Theorem 1) that the symmetry group of the linearized system $`\ddot{𝜹}=\mathrm{J}(t)𝜹`$ turns out to be precisely the symmetry group $`G`$ of the considered bush B$`[G]`$. This result allows one to apply the well-known Wigner theorem about the specific structure of the matrix ($`\mathrm{J}(t)`$, in our case) commuting with all the matrices of a fixed representation (mechanical representation, in our case) of a given group. According to the above theorem one can split effectively the linearized system $`\ddot{𝜹}=\mathrm{J}(t)𝜹`$ into a number of independent subsystems of differential equations with time-dependent coefficients.
We want to emphasize that this symmetry-related method for splitting the linearized systems arising in the linear stability analysis of the dynamical regimes is suitable for *arbitrary* nonlinear mechanical systems with discrete symmetry. Such a decomposition (splitting) of the linearized system $`\ddot{𝜹}=\mathrm{J}(t)𝜹`$ is especially important for the *multidimensional* bushes of modes, describing *quasiperiodic* vibrational regimes, which cannot be treated with the aid of the Floquet method. Indeed, in this case, we need to integrate the differential equations with time-dependent coefficients over large time interval, unlike the case of periodic regimes where we can solve the appropriate differential equations over only one period to construct the monodromy matrix.
The above method is applied for studying the stability of some dynamical regimes (bushes of modes) in the monoatomic chains. For this specific mechanical systems, we prove Theorem 2 which allows one to find very simply the upper bound of dimensions of the independent subsystems obtained after splitting the linearized system $`\ddot{𝜹}=\mathrm{J}(t)𝜹`$. Indeed, according to this theorem, the dimension of each such subsystem does not exceed the integer $`m`$ determining the ratio of the volumes of the primitive cell of the chain in the vibrational state, corresponding to the given bush B$`[G]=`$B$`[\widehat{a}^m,\mathrm{}]`$, and the equilibrium state.
Taking into account any other symmetry elements of the considered bush allows to reduce the dimensions of at least some of the above discussed subsystems. We illustrate this fact comparing the stability analysis of the $`\pi `$-mode (zone boundary mode) for the FPU-$`\alpha `$ and FPU-$`\beta `$ chains.
###### Acknowledgements.
We are very grateful to Prof. V.P. Sakhnenko for useful discussions and for his friendly support, and to O.E. Evnin for his valuable help with the language corrections in the text of this paper.
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# Holographic Duals of a Family of 𝒩=1^∗ Fixed Points
## 1 Introduction
Motivated by the AdS/CFT duality , there has been considerable interest in finding explicit supergravity solutions that correspond to conformal field theories with N=1 supersymmetry in four dimensions. In this paper we solve a long standing problem in this context which was originally posed in : We find the conjectured family of solutions that correspond to infra-red fixed-points that interpolate between the solution of Pilch and Warner , and the solution of Romans that is the basis of what has become known as the Klebanov-Witten (KW) point .
To be more precise, the Klebanov and Witten argued that if one starts with the $`𝒩=2`$, four-dimensional, $`\widehat{A}_1`$ quiver gauge theory and breaks it to an $`𝒩=1`$ supersymmetric field theory by introducing a (unique) $`SO(4)`$ invariant superpotential, the theory will flow to an $`𝒩=1`$ superconformal fixed point in the infra-red and this fixed point is dual to the solution of IIB supergravity on $`\mathrm{AdS}_5\times T^{1,1}`$ . Similarly, it was argued in that the same $`\widehat{A}_1`$ quiver gauge theory would, under a particular $`SO(3)`$ invariant superpotential, flow to another $`𝒩=1`$ superconformal fixed point whose supergravity dual is the ($`_2`$ orbifold of) the Pilch-Warner (PW) solution whose existence was first discovered via five-dimensional supergravity .
More generally, it was argued in , using the non-perturbative methods of Leigh and Strassler , that there is a family of four-dimensional $`𝒩=1`$ superconformal field theories (SCFT’s) that continuously interpolate between the KW flow and the PW flow, and that this family preserves at least an $`SO(3)`$ global symmetry. Indeed, also investigated the corresponding five-dimensional gauged supergravity solutions that were expected to capture the relevant sectors of the IIB supergravity dual of the family of flows. From the five-dimensional perspective, the existence of the family of flows and of the family of IR fixed points was almost a triviality. There was, however, an important caveat: There are no consistent truncation theorems for this more general class of five-dimensional supergravity theories, and so the five-dimensional result were very suggestive, but did not prove that there had to be corresponding ten-dimensional solutions. The search for this family of solutions within IIB supergravity has been rather long and surprisingly difficult, and here we will prove that family exists by reducing the problem to a system of ordinary differential equations and exhibiting numerical solutions.
Much of the technology for finding supersymmetric solutions to supergravities in various dimensions relies on solving the supersymmetry variations and the Bianchi identities, one can a postieri check that the field equations are satisfied. A general formalism for analyzing the Killing spinor equations is that of $`G`$-structures. For IIB supergravity, this works extremely well when the internal manifold has $`SU(3)`$ structure but for backgrounds with only $`SU(2)`$ structure (which is the structure appearing in the current work) that methodology is too cumbersome at present . A more pragmatic approach developed by two of the current authors and their collaborators , is to use the physics of the problem to make an Ansatz for the Killing spinors as well as the metric and form fields. We will follow the latter approach and find that the symmetries of the problem sufficiently restrict the form of the Ansatz such that the full solution can be obtained. More specifically, the“supersymmetry bundle” is a four-dimensional subspace of the of the 32 real components of the spinors, and we can use the symmetries and a specific combination of the gravitino variation equations to define an eight-dimensional subspace that contains the Killing spinors. We then parametrize the supersymmetries within this eight-dimensional subspace in a manner that is equivalent to the dielectric deformation of the canonical $`D3`$-brane projector . Having found the supersymmetries, one can then build the rest of the solution from the Killing spinor equation.
The solution of IIB supergravity on $`\mathrm{AdS}_5\times T^{1,1}`$ is a Freund-Rubin Ansatz with constant dilaton-axion and vanishing three-form flux. One can re-cast this solution in terms of $`D3`$ branes on the conifold, and the metric transverse to the branes is thus Kähler and Ricci flat. It therefore possesses a rather trivial $`SU(3)`$ structure. The PW solution is a warped Freund-Rubin Ansatz with constant dilaton-axion and non-vanishing three-form flux. The PW metric is neither Ricci flat nor Kähler but it is equipped with an integrable complex structure, namely that of $`A_1\times `$ . It has two globally-defined spinors and as such has only $`SU(2)`$ structure. We find that the interpolating solution also has only an $`SU(2)`$ structure. The surprise is that even though the two end points of our interpolation have a trivial dilaton-axion, the interpolating solutions themselves have a non-trivial dilaton-axion. It also seems that the interpolating family lacks a integrable complex structure.
It is worth mentioning the interesting recent work in which the authors use the eight-dimensional duality group to generate new solutions that can be easily lifted to ten dimensions. For the supergravity duals to SCFT’s they are able to identify the exactly marginal operator in the field theory, thus providing a holographic check of the methods of Leigh and Strassler. Our scenario falls out of the scope of the powerful methods employed there since it lacks the required two $`U(1)`$ non-$`R`$ symmetries.
This paper is organized as follows: In section 2, we review the relevant field theory, and in particular discuss the symmetries. The symmetries of the supergravity background are discussed in Section 3. In Section 4, we reduce the problem to five dimensions, enforcing the $`\mathrm{AdS}_5`$ factor in the ten-dimensional background. Section 5 contains a review of the KW and PW solutions. Sections 6 and 7 contain the main calculations: We derive the BPS equations from the most general Ansatz which preserves the relevant symmetries and reduce this system to three first order, non-linear ordinary differential equations. In section 7 we establish that there is indeed a one parameter family of regular solutions to these BPS equations and solve them numerically. We indeed show that they interpolate between the KW and PW solutions. Those who are interested in the main result should therefore jump to sections 6 and 7. Section 8 contains a discussion of the central charge of each gauge theory in the family from the perspective of the dual supergravity theory. We show analytically that the central charge has the correct constant value across the entire family of solutions. Finally, there are several appendices containing spinor conventions and computational details.
## 2 Field theory considerations
The conformal field theory we are considering in this paper is a non-trivial IR fixed point of a mass deformed $`𝒩=2`$ quiver gauge theory . The UV field theory has an $`SU(N)\times SU(N)`$ gauge group two bi-fundamental hypermultiplets, one in the $`(N,\overline{N})`$ and one in the $`(\overline{N},N)`$. In $`𝒩=1`$ language the first hypermultiplet decomposes into two chiral multiplets $`(A_1,B_1)`$ and the second hypermultiplet decomposes into two chiral multiplets $`(A_2,B_2)`$.
The superpotential of this theory is
$$W=\mathrm{Tr}\left(\varphi _1(A_1B_1B_2A_2)\right)+\mathrm{Tr}\left(\varphi _2(A_2B_2B_1A_1)\right).$$
(2.1)
This theory has an $`SU(2)\times SU(2)_R\times U(1)_R`$ continous global symmetry. The two hypermultiplets form a doublet under the $`SU(2)`$ flavor symmetry.
This theory can be deformed by mass terms for the adjoint scalars
$$\mathrm{\Delta }W=\frac{m_1}{2}\mathrm{Tr}\left(\varphi _1^2\right)+\frac{m_2}{2}\mathrm{Tr}\left(\varphi _2^2\right).$$
(2.2)
This deformation breaks the continuous global symmetry to $`SU(2)\times U(1)_R`$. The $`U(1)_R`$ symmetry is actually a combination of the $`U(1)_R`$ symmetry and a $`U(1)`$ subgroup of the $`SU(2)_R`$ symmetry of the $`𝒩=2`$ theory. The R-charges of the fields are
$$\begin{array}{ccc}\varphi _i& A_i& B_i\\ & & \\ 1& \frac{1}{2}& \frac{1}{2}\end{array}$$
(2.3)
This field theory also has a $`_4`$ discrete symmetry, with a generator $`\omega `$ which acts as a charge conjugation
$`\varphi _i`$ $``$ $`\varphi _i^t`$ (2.4)
$`A_i`$ $``$ $`iA_{i+1}^t`$ (2.5)
$`B_i`$ $``$ $`iB_{i+1}^t`$ (2.6)
It is easy to see that the $`_4`$ symmetry commutes with the continous symmetries and so the global symmetry of the theory is $`SU(2)\times _4\times U(1)_R`$. However, $`\omega ^2`$ and the center of $`SU(2)`$ simply negate $`A_i`$ and $`B_i`$, and in the supergravity dual we will consider only gauge-invariant bilinears of the fields $`A,B`$. Thus these generators will act trivially in supergravity which means that the symmetry of the supergravity theory<sup>1</sup><sup>1</sup>1Indeed, even within the $`SU(N)`$ gauge theory, for $`N`$ even, negating $`A_i`$ and $`B_i`$ is in the center of the $`SU(N)`$ gauge groups and so the symmetry of (perturbative) physical states of the field theory will also be $`SO(3)\times _2\times U(1)_R`$. will be $`SO(3)\times _2\times U(1)_R`$.
Below the mass scale given by $`m_1`$ and $`m_2`$ one can integrate out the adjoint scalars $`\varphi _1`$ and $`\varphi _2`$ and the low energy superpotential is given by
$$W=\lambda _1\mathrm{Tr}\left((A_1B_1B_2A_2)^2\right)+\lambda _2\mathrm{Tr}\left((A_2B_2B_1A_1)^2\right).$$
(2.7)
The low energy effective action has the two gauge couplings $`\tau _i`$ and the two quartic superpotential couplings $`\lambda _i`$.
The deformed theory is believed to flow to a non-trivial IR fixed point. Vanishing of the $`\beta `$-functions for all the couplings requires
$$\gamma _{A_i}(\tau _1,\tau _2,\lambda _1,\lambda _2)+\gamma _{B_i}(\tau _1,\tau _2,\lambda _1,\lambda _2)+\frac{1}{2}=0.$$
(2.8)
This is two equations for four unknowns. However, the $`SU(2)\times _4`$ symmetry implies that the functional form of all the anomalous dimensions is the same
$$\gamma =\gamma _{A_i}=\gamma _{A_{i+1}}=\gamma _{B_i}=\gamma _{B_{i+1}}.$$
(2.9)
From this we conclude that the vanishing of the $`\beta `$-function implies only one constraint
$$\gamma +\frac{1}{4}=0$$
(2.10)
for four unknowns. We expect the moduli space of IR theories to have three complex dimensions.
The central charges for such theories have been calculated in . The ratio of the central charges of the IR theory and the UV theory is
$$\frac{c_{(IR)}}{c_{(UV)}}=\frac{27}{32}.$$
(2.11)
## 3 Realizing the global symmetries within supergravity
In the following we want to construct supergravity backgrounds that are holographic duals to field theories with a given global symmetry algebra. The global symmetries have to be realized as symmetries of the background and this leads to powerful constraints on the background. One set of constraints comes from the existence of the symmetry generators. The other set of constraints comes from the commutation relations of the symmetry generators.
Type IIB supergravity has six different gauge symmetries. General coordinate transformations which we restrict to the isometries generated by Killing vectors, $`\delta (\xi )`$; local Lorentz transformations, $`\delta (l)`$; the U(1) R-symmetry, $`\delta (\mathrm{\Sigma })`$; the gauge transformations of the two-form and four-form potentials, $`\delta (\mathrm{\Lambda }^{(1)})`$ and $`\delta (\mathrm{\Lambda }^{(3)})`$, and the supersymmetry transformation $`\delta (ϵ)`$. There is also a global $`SU(1,1)`$ symmetry, which in string theory is broken to $`SL(2,)`$.
A background has a global symmetry generated by some specific symmetry generators, $`(\xi ,l,\mathrm{\Sigma },\mathrm{\Lambda }^{(1)},\mathrm{\Lambda }^{(3)},ϵ)`$, provided that this transformation leaves the background invariant<sup>2</sup><sup>2</sup>2There are also the $`SL(2,)`$ actions, but those are discrete symmetries.. Global supersymmetries have to be generated just by an $`ϵ`$ and global bosonic symmetries will be generated by a combination $`(\xi ,l,\mathrm{\Sigma },\mathrm{\Lambda }^{(1)},\mathrm{\Lambda }^{(3)})`$.
### 3.1 Continous bosonic symmetries
The non-trivial, physical bosonic symmetries of the background must involve a transformation by an isometry, or a Killing vector<sup>3</sup><sup>3</sup>3One can see this from the fact, that a transformation generated by $`(\xi =0,l,\mathrm{\Sigma },\mathrm{\Lambda }^{(1)},\mathrm{\Lambda }^{(3)})`$ will not leave any field configuration invariant., $`\xi `$. The vielbein only transforms under both general coordinate and local Lorentz transformations and its invariance typically requires a compensating local Lorentz transformation that depends on the choice of the vielbein. For this reason it is often useful, wherever possible, to choose a vielbein made of invariant one-forms.
The coset fields $`V_\pm ^\alpha `$, which describe the dilaton and axion, transform under the global $`SU(1,1)`$ symmetry and locally under general coordinate transformations and the $`U(1)`$ R-symmetry of the IIB theory. The invariance of the coset fields requires
$$(_\xi \pm i\mathrm{\Sigma })V_\pm ^\alpha =0.$$
(3.1)
From this it is easy to derive that the gradient of the dilaton-axion field in the $`\xi `$ direction is vanishing
$$P_\xi =0$$
(3.2)
and that $`\mathrm{\Sigma }`$ is the $`\xi `$ component $`Q_\xi `$ of the $`U(1)`$ connection. Equation (3.1) also implies, that at least one of the $`V_\pm ^\alpha `$ is a non-vanishing section of the associated line bundle over an orbit of $`\xi `$. This implies, that one can choose a trivialization of the $`U(1)`$ bundle over an orbit of $`\xi `$ and thereby render the connection, $`Q_\xi `$, trivial:
$$Q_\xi =0.$$
(3.3)
With this choice of trivialization the action of the symmetry on the field strengths $`G^{(3)}`$ and $`F^{(5)}`$ is just the Lie derivative. For this reason $`G^{(3)}`$ and $`F^{(5)}`$ have to be invariant forms
$$\mathrm{\pounds }_\xi G^{(3)}=0\mathrm{and}\mathrm{\pounds }_\xi F^{(5)}=0.$$
(3.4)
We will not discuss the gauge transformations $`\delta (\mathrm{\Lambda }^{(1)},\mathrm{\Lambda }^{(3)})`$ here because the supersymmetry variations, the Bianchi identities and the equations of motion depend only on the field strengths $`G^{(3)}`$ and $`F^{(5)}`$.
From the above discussion it follows, that the Killing vectors have to satisfy the bosonic Lie algebra of the global symmetry group of the background
$$[\delta (u_1),\delta (u_2)]=\delta ([u_1,u_2]).$$
(3.5)
This implies that the background is a fibration of a product of coset spaces and group manifolds over a possibly non-trivial base.
### 3.2 Supersymmetries
The supersymmetries are generated by Killing spinors $`ϵ`$. In a purely bosonic background the requirement of the existence of a global supersymmetry is the vanishing of the dilatino and the gravitino variation.
Before looking at the dilatino and gravitino variation it is useful to look at the commutators of the supersymmetry generators with other symmetry generators.
$$[\delta (g),\delta (ϵ)]=[\delta (\xi ),\delta (ϵ)]+[\delta (l),\delta (ϵ)]+[\delta (\mathrm{\Sigma }),\delta (ϵ)]=\delta \left(\left(_\xi +\frac{1}{4}l_{rs}\gamma ^{rs}\frac{i}{2}Q_\xi \right)ϵ\right),$$
(3.6)
where $`l`$ is a “Lie connection.” If the vielbein is given in terms of invariant forms and the U(1) connection is chosen trivially, then the above expression reduces to the ordinary derivative. For later convenience we define the Lie derivative of $`ϵ`$ by this derivative operator:
$$\mathrm{\pounds }_\xi ϵ=\left(_\xi +\frac{1}{4}l_{rs}\gamma ^{rs}\frac{i}{2}Q_\xi \right)ϵ.$$
(3.7)
This gives rise to the differential equation
$$\mathrm{\pounds }_\xi ϵ=gϵ$$
(3.8)
This allows to determine the dependence of $`ϵ`$ on the directions given by the symmetries.
There are also powerful constraints coming from the anti-commutator of two supersymmetries
$$\{\delta (ϵ_1),\delta (ϵ_2)\}=\delta (\{ϵ_1,ϵ_2\}).$$
(3.9)
This implies, that
$$\xi ^\mu =2\mathrm{I}\mathrm{m}(\overline{ϵ}_1\gamma ^\mu ϵ_2)$$
(3.10)
is the Killing vector associated to $`\{ϵ_1,ϵ_2\}`$ and that
$$l^{rs}=\omega _\xi {}_{}{}^{rs}\frac{1}{3}F^{rsmnp}\mathrm{Re}(\overline{ϵ}_1\gamma _{mnp}ϵ_2)+\frac{3}{4}\mathrm{Im}(G^{rsm}\overline{ϵ}_1\gamma _mϵ_2^{}+\frac{1}{18}G_{mnp}\overline{ϵ}_1\gamma ^{rsmnp}ϵ_2^{})$$
(3.11)
is the local Lorentz transformation associated to $`\{ϵ_1,ϵ_2\}`$. The relation including $`\mathrm{\Sigma }`$ is trivially satisfied.
### 3.3 Discrete symmetries
Discrete symmetries can be composed out of global diffeomorphisms, local Lorentz transformations, gauge transformations for the form fields, $`U(1)_R`$ symmetry transformations and $`SL(2,)`$ transformations. The commutation relations with the other global symmetries are again given by the field theory.
The constraints from discrete symmetries are especially powerful when the global diffeomorphism leaves the orbits of the continous symmetries invariant. If this is the case, the discrete symmetry implies powerful projection conditions on the fields and supersymmetry generators. We will see an explicit example of this below.
## 4 Reduction to a five-dimensional problem
### 4.1 Decomposing the metric and spinors
The bosonic part of the four-dimensional, $`𝒩=1`$ superconformal algebra is $`SO(2,4)\times U(1)`$. This bosonic symmetry is realized by Killing vectors in the ten-dimensional geometry. The geometry is then $`AdS_5`$, which is covering space of $`SO(2,4)/SO(1,4)`$, warped over an internal five-manifold. The internal five-manifold itself is a $`S^1`$ fibration over a four-manifold, $`X_4`$. The $`S^1`$ must be a Killing direction dual to the $`R`$-symmetry action.
We adopt the following index conventions: Capital Latin letters denote ten-dimensional indices ($`0,\mathrm{},9`$), small Greek letters denote the five-dimesional indices in the $`AdS^5`$ ($`0,\mathrm{},4`$) and small Latin letters denote the internal indices ($`5,\mathrm{},9`$). A hat denotes ten-dimensional frame indices, a tilde denotes five-dimensional frame indices in $`AdS_5`$ and a check denotes five-dimensional frame indices in the internal space. The warped $`AdS_5`$ leads to a vielbein Ansatz of the form:
$`e^{\widehat{\mu }}`$ $`=`$ $`\mathrm{\Omega }e_{(e)}^{\stackrel{~}{\mu }}\mathrm{for}\mu =0,\mathrm{},4`$ (4.1)
$`e^{\widehat{m}}`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Omega }}}e_{(i)}^{\stackrel{ˇ}{m}}\mathrm{for}m=5,\mathrm{},9`$ (4.2)
where $`e_{(e)}^{\stackrel{~}{\mu }}`$ is a vielbein for $`AdS_5`$ of unit curvature radius and $`e_{(i)}^{\stackrel{ˇ}{m}}`$ is a vielbein for the internal manifold.
The spinors of IIB supergravity must similarly decompose into spinors on $`AdS_5`$ and on the internal five-manifold. We will analyze this in detail, and we need to recall some basic facts about spinors in various dimensions. More information may be found in Appendix A.
Recall that in the IIB theory one can impose a Majorana-Weyl condition on a spinor to reduce it to 16 real components. It is most convenient to represent the 32 components of the $`𝒩=2`$ supersymmetry of the IIB theory in terms of a complex Weyl spinor. Our task will be to decompose this into components along the two five-manifolds. To do this it will be important to recall how complex conjugation acts on spinors. Given a set of $`\gamma `$-matrices, complex conjugation maps them into an equivalent set, and so there is a matrix, $`B`$, that will generically conjugate the $`\gamma _A^{}`$ back to the $`\gamma _A`$. By the same token, to map a spinor, $`\mathrm{\Psi }`$, to its complex conjugate representation one must accompany the conjugation by the action of $`B`$. Thus, the conjugate spinor, $`\mathrm{\Psi }^{}`$, is defined by:
$$\mathrm{\Psi }^{}B^1\mathrm{\Psi }^{}.$$
(4.3)
The form and properties of $`B`$ depend upon the dimension and signature of the metric and upon $`\gamma `$-matrix conventions. In the IIB theory one can adopt conventions in which $`B`$ is the identity matrix (as in )), but we will keep our expressions convention independent and adopt the notation (4.3). In five Lorentzian dimensions, $`B`$ is necessarily non-trivial and may be thought of as a symplectic form. Indeed, this fact lies at the heart of the symplectic Majorana condition of five-dimensional supersymmetric theories.
In the following we will use the the notation (4.3) to denote the conjugate spinor in all dimensions and metric signatures.
#### 4.1.1 Killing spinors on $`AdS_5`$
The ten-dimensional Killing spinor is a complex, chiral spinor ($`\gamma _{(10)}ϵ=ϵ`$). Since it has to respect the symmetries of $`AdS_5`$, it has to be built out of five-dimensional Killing spinors. The five-dimensional Killing spinor equation is:
$$(D_\mu \pm \frac{i}{2}\gamma _\mu )\zeta =0,$$
(4.4)
for either choice of sign. The distinct signs determine the transformation properties under the conformal group, $`SO(2,4)`$. That is, solutions with a plus (respectively, minus) sign transform in the $`4`$ (respectively, $`\overline{4}`$) of $`SO(2,4)`$. If $`\zeta `$ is any $`SO(1,4)`$ four-spinor satisfying (4.4) for one choice of sign, it is easy to see that $`\zeta ^{}`$ is a solution to (4.4) with the opposite sign.
One can also check that $`2\mathrm{R}\mathrm{e}(\overline{\zeta }_1\gamma ^\mu \zeta _2)`$ are Killing vector fields generating the $`AdS_5`$ group and that $`\mathrm{Re}(\overline{\zeta }_1\zeta _2)`$ generates the $`U(1)_R`$ symmetry in accordance with the four-dimensional, $`𝒩=1`$ superconformal algebra. This is because the superconformal algebra implies that the bosonic symmetry generators appear in the $`4\overline{4}=115`$. On the other hand, expressions like $`\mathrm{Re}(\overline{\zeta }_1\zeta _2^{})`$ and $`\mathrm{Re}(\stackrel{~}{\zeta }_1\gamma ^\mu \zeta _2)`$ are not related Killing vectors or other bosonic symmetry generators.
#### 4.1.2 The ten-dimensional Killing spinors
The ten-dimensional Killing spinors can be decomposed as
$$ϵ_\zeta =\mathrm{\Omega }^{\frac{1}{2}}\left(\begin{array}{c}\zeta \chi ^{(1)}+(\zeta ^{})(\chi ^{(2)}{}_{}{}^{})\\ 0\end{array}\right),$$
(4.5)
where $`\zeta `$ is a Killing spinor in $`AdS_5`$ which does not depend on the internal coordinates and $`\chi ^{(i)}`$ are independent internal five-dimensional spinors which only depend on the internal coordinates.
We can now compute the Killing vectors $`\mathrm{Re}(\overline{ϵ}_1\gamma ^Mϵ_2)`$
$`\mathrm{Re}(\overline{ϵ}_1\gamma ^\mu ϵ_2)`$ $`=`$ $`\mathrm{\Omega }\mathrm{Re}(\overline{\zeta }_1\gamma ^\mu \zeta _2)(\overline{\chi }^{(1)}\chi ^{(1)}+\overline{\chi }^{(2)}\chi ^{(2)}),`$ (4.6)
$`\mathrm{Re}(\overline{ϵ}_1\gamma ^mϵ_2)`$ $`=`$ $`\mathrm{\Omega }\mathrm{Re}(\overline{\zeta }_1\zeta _2)(\overline{\chi }^{(1)}\gamma ^m\chi ^{(1)}+\overline{\chi }^{(2)}\gamma ^m\chi ^{(2)}).`$ (4.7)
It is interesting to note that cross terms like $`\stackrel{~}{\chi }^{(1)}\chi ^{(2)}`$ cancel out in this expression. The foregoing equations also give rise to normalization conditions for the $`\chi ^{(i)}`$. The condition coming from the normalization of the Killing vectors parallel to the $`AdS_5`$ is:
$$\overline{\chi }^{(1)}\chi ^{(1)}+\overline{\chi }^{(2)}\chi ^{(2)}=1.$$
(4.8)
Similarly, the Killing vector of the form (4.7) along the internal manifold must be that of the $`U(1)_R`$ symmetry, and so we must have:
$$\frac{3}{2}\frac{}{\varphi }=\mathrm{\Omega }(\overline{\chi }^{(1)}\gamma ^m\chi ^{(1)}+\overline{\chi }^{(2)}\gamma ^m\chi ^{(2)})e_{\widehat{m}},$$
(4.9)
where $`\varphi `$ is an internal coordinate.
Finally, we can determine the $`\varphi `$ dependence of the internal spinors $`\chi ^{(i)}`$. Since the $`\varphi `$ direction realizes the $`U(1)_R`$ symmetry, we have to impose
$$\mathrm{\pounds }_\frac{}{\varphi }ϵ_\zeta =ϵ_{i\zeta },$$
(4.10)
which is equivalent to the five-dimensional spinor $`\zeta `$ having charge 1 under the $`U(1)_R`$ symmetry. This leads to
$$\mathrm{\pounds }_\frac{}{\varphi }\chi ^{(i)}=i\chi ^{(i)}.$$
(4.11)
### 4.2 The dilatino variation
The dilatino variation is given by
$$\delta \lambda =iP_{\widehat{M}}\gamma ^Mϵ^{}\frac{i}{24}G_{\widehat{M}\widehat{N}\widehat{P}}\gamma ^{MNP}ϵ.$$
(4.12)
Poincaré invariance requires that $`P`$ and $`G`$ only have components in the internal directions. This leads to the equation:
$$0=P_{\widehat{m}}\gamma ^mϵ^{}\frac{1}{24}G_{\widehat{m}\widehat{n}\widehat{p}}\gamma ^{mnp}ϵ.$$
(4.13)
Inserting the form of the Killing spinor (4.5) and realizing that $`\zeta `$ and $`\zeta ^{}`$ may be considered as independent variables, we get the two five-dimensional equations:
$`\delta \lambda ^{(1)}=P_{\stackrel{ˇ}{m}}\gamma _{(i)}^m\chi ^{(2)}+{\displaystyle \frac{\mathrm{\Omega }^2}{24}}G_{\stackrel{ˇ}{m}\stackrel{ˇ}{n}\stackrel{ˇ}{p}}\gamma _{(i)}^{mnp}\chi ^{(1)}=0,`$ (4.14)
$`\delta \lambda ^{(2)}=P_{\stackrel{ˇ}{m}}\gamma _{(i)}^m\chi ^{(1)}{}_{}{}^{}+{\displaystyle \frac{\mathrm{\Omega }^2}{24}}G_{\stackrel{ˇ}{m}\stackrel{ˇ}{n}\stackrel{ˇ}{p}}\gamma _{(i)}^{mnp}\chi ^{(2)}{}_{}{}^{}=0.`$ (4.15)
Since the background fields are independent of the $`U(1)_R`$ direction, these equations reduce to spinor equations on the four-dimensional base, $`X_4`$, of the $`S^1`$ fibration that makes up the internal manifold. It is also easy to show that the component of $`P_{\stackrel{ˇ}{m}}`$ along the $`U(1)_R`$ Killing vector must vanish. This result is expected from the $`U(1)_R`$ invariance, but can be deduced explicitly from (4.14) and (4.15) as follows: Multiply the first equation by $`\overline{\chi }^{(2)}`$, transpose the second equation and multiply it by $`\chi ^{(1)}`$ and add the two.
### 4.3 The gravitino variation
The gravitino variation is :
$$\delta \psi _{\widehat{M}}=D_{\widehat{M}}ϵ+\frac{i}{480}F_{\widehat{P}\widehat{Q}\widehat{R}\widehat{S}\widehat{T}}\gamma ^{PQRST}\gamma _Mϵ+\frac{1}{96}G_{\widehat{P}\widehat{Q}\widehat{R}}\left(\gamma _M{}_{}{}^{PQR}9\delta _M^P\gamma ^{QR}\right)ϵ^{},$$
(4.16)
where the covariant derivative is given by
$$D_{\widehat{M}}ϵ=_{\widehat{M}}ϵ+\frac{1}{4}\omega _{\widehat{M}\widehat{P}\widehat{Q}}\gamma ^{PQ}ϵ\frac{i}{2}Q_{\widehat{M}}ϵ.$$
(4.17)
In order to continue, we need to determine the ten-dimensional spin connection in terms of the warp factor $`\mathrm{\Omega }`$ and the five-dimensional spin connection:
$`\omega _{\widehat{\mu }\widehat{\nu }\widehat{\rho }}`$ $`=`$ $`\mathrm{\Omega }^1\omega _{(e)}{}_{\stackrel{~}{\mu }\stackrel{~}{\nu }\stackrel{~}{\rho }}{}^{},`$ (4.18)
$`\omega _{\widehat{\mu }\widehat{\nu }\widehat{r}}`$ $`=`$ $`_{\stackrel{ˇ}{r}}\mathrm{\Omega }\eta _{\widehat{\mu }\widehat{\nu }},`$ (4.19)
$`\omega _{\widehat{\mu }\widehat{n}\widehat{r}}`$ $`=`$ $`0,`$ (4.20)
$`\omega _{\widehat{m}\widehat{\nu }\widehat{\rho }}`$ $`=`$ $`0,`$ (4.21)
$`\omega _{\widehat{m}\widehat{\nu }\widehat{r}}`$ $`=`$ $`0,`$ (4.22)
$`\omega _{\widehat{m}\widehat{n}\widehat{r}}`$ $`=`$ $`\mathrm{\Omega }\omega _{(i)}{}_{\stackrel{ˇ}{m}\stackrel{ˇ}{n}\stackrel{ˇ}{r}}{}^{}_{\stackrel{ˇ}{n}}\mathrm{\Omega }\delta _{\stackrel{ˇ}{m}\stackrel{ˇ}{r}}+_{\stackrel{ˇ}{r}}\mathrm{\Omega }\delta _{\stackrel{ˇ}{m}\stackrel{ˇ}{n}},`$ (4.23)
where $`\omega _{(e)}`$ is the spin connection on $`AdS_5`$ and $`\omega _{(i)}`$ is the spin connection on the internal manifold. Also note that
$$Q_{\widehat{\mu }}=0.$$
(4.24)
The self dual five-form flux can be written as
$$F^{(5)}=fe^{\widehat{0}\mathrm{}\widehat{4}}+fe^{\widehat{5}\mathrm{}\widehat{9}},$$
(4.25)
where $`f`$ only depends on the internal coordinates. The Bianchi identity for $`F^{(5)}`$ reduces, for such a compactification, to
$$dF^{(5)}=0,$$
(4.26)
which implies
$$f=\frac{f_0}{\mathrm{\Omega }^5},$$
(4.27)
where $`f_0`$ is an integration constant.
Now we can determine the gravitino variations with $`M=0,\mathrm{},4`$ a similar argument as for the dilatino variation leads to
$`\delta \psi _0^{(1)}={\displaystyle \frac{i}{2\mathrm{\Omega }^2}}\chi ^{(1)}+{\displaystyle \frac{1}{2}}_{\stackrel{ˇ}{r}}\mathrm{log}\mathrm{\Omega }\gamma _{(i)}^r\chi ^{(1)}+{\displaystyle \frac{if_0}{2\mathrm{\Omega }^6}}\chi ^{(1)}{\displaystyle \frac{\mathrm{\Omega }^2}{96}}G_{\stackrel{ˇ}{p}\stackrel{ˇ}{q}\stackrel{ˇ}{r}}\gamma _{(i)}^{pqr}\chi ^{(2)}=0,`$ (4.28)
$`\delta \psi _0^{(2)}={\displaystyle \frac{i}{2\mathrm{\Omega }^2}}\chi ^{(2)}{}_{}{}^{}+{\displaystyle \frac{1}{2}}_{\stackrel{ˇ}{r}}\mathrm{log}\mathrm{\Omega }\gamma _{(i)}^r\chi ^{(2)}{}_{}{}^{}+{\displaystyle \frac{if_0}{2\mathrm{\Omega }^6}}\chi ^{(2)}{}_{}{}^{}{\displaystyle \frac{\mathrm{\Omega }^2}{96}}G_{\stackrel{ˇ}{p}\stackrel{ˇ}{q}\stackrel{ˇ}{r}}\gamma _{(i)}^{pqr}\chi ^{(1)}{}_{}{}^{}=0.`$ (4.29)
Similarly, the gravitino variations with $`M=5,\mathrm{},10`$ lead to
$`\delta \psi _{\stackrel{ˇ}{m}}^{(1)}=D_{\stackrel{ˇ}{m}}\chi ^{(1)}+{\displaystyle \frac{1}{2}}_{\stackrel{ˇ}{m}}\mathrm{log}\mathrm{\Omega }\chi ^{(1)}{\displaystyle \frac{1}{2}}_{\stackrel{ˇ}{r}}\mathrm{log}\mathrm{\Omega }\gamma _{(i)}{}_{m}{}^{}{}_{}{}^{r}\chi _{}^{(1)}`$
$`{\displaystyle \frac{if_0}{2\mathrm{\Omega }^6}}\gamma _{(i)}{}_{m}{}^{}\chi _{}^{(1)}{\displaystyle \frac{\mathrm{\Omega }^2}{96}}G_{\stackrel{ˇ}{p}\stackrel{ˇ}{q}\stackrel{ˇ}{r}}\gamma _{(i)}{}_{m}{}^{}{}_{}{}^{pqr}\chi _{}^{(2)}+{\displaystyle \frac{3\mathrm{\Omega }^2}{32}}G_{\stackrel{ˇ}{m}\stackrel{ˇ}{q}\stackrel{ˇ}{r}}\gamma _{(i)}{}_{}{}^{qr}\chi _{}^{(2)}`$ $`=`$ $`0`$ (4.30)
$`\delta \psi _{\stackrel{ˇ}{m}}^{(2)}=D_{\stackrel{ˇ}{m}}\chi ^{(2)}{}_{}{}^{}+{\displaystyle \frac{1}{2}}_{\stackrel{ˇ}{m}}\mathrm{log}\mathrm{\Omega }\chi ^{(2)}{}_{}{}^{}{\displaystyle \frac{1}{2}}_{\stackrel{ˇ}{r}}\mathrm{log}\mathrm{\Omega }\gamma _{(i)}{}_{m}{}^{}{}_{}{}^{r}\chi _{}^{(2)}{}_{}{}^{}`$
$`{\displaystyle \frac{if_0}{2\mathrm{\Omega }^6}}\gamma _{(i)}{}_{m}{}^{}\chi _{}^{(2)}{}_{}{}^{}{\displaystyle \frac{\mathrm{\Omega }^2}{96}}G_{\stackrel{ˇ}{p}\stackrel{ˇ}{q}\stackrel{ˇ}{r}}\gamma _{(i)}{}_{m}{}^{}{}_{}{}^{pqr}\chi _{}^{(1)}{}_{}{}^{}+{\displaystyle \frac{3\mathrm{\Omega }^2}{32}}G_{\stackrel{ˇ}{m}\stackrel{ˇ}{q}\stackrel{ˇ}{r}}\gamma _{(i)}{}_{}{}^{qr}\chi _{}^{(1)}^{}`$ $`=`$ $`0`$ (4.31)
## 5 Known solutions
### 5.1 The $`T^{1,1}`$ solution
The $`T^{1,1}`$ space is the intersection of the conifold $`z_1^2+z_2^2+z_3^2+z_4^2=0`$ with the unit sphere. This can be obtained by applying $`SO(3)\times U(1)`$ transformations on vectors of the form
$$(z_1^{(0)},\mathrm{},z_4^{(0)})=(1,i\mathrm{cos}\theta ,0,i\mathrm{sin}\theta ).$$
(5.1)
We can use the $`SO(3)\times U(1)`$ rotations to ensure that $`\mathrm{cos}\theta ,\mathrm{sin}\theta 0`$, and so the manifold is covered if one takes $`0\theta \frac{\pi }{2}`$. Applying infinitesimal transformations
$$\left(\begin{array}{cccc}1+id\varphi & \sigma ^3& \sigma ^2& 0\\ \sigma _3& 1+id\varphi & \sigma ^1& 0\\ \sigma ^2& \sigma ^1& 1+id\varphi & 0\\ 0& & 0& 1+id\varphi \end{array}\right)$$
(5.2)
leads to
$$d\stackrel{}{z}=\left(\begin{array}{c}dz_1\\ dz_2\\ dz_3\\ dz_4\end{array}\right)=\left(\begin{array}{c}id\varphi i\mathrm{cos}\theta \sigma ^3\\ \sigma ^3\mathrm{cos}\theta d\varphi i\mathrm{sin}\theta d\theta \\ \sigma ^2+i\mathrm{cos}\theta \sigma ^1\\ \mathrm{sin}\theta d\varphi +i\mathrm{cos}\theta d\theta \end{array}\right).$$
(5.3)
The metric then takes the form
$$ds^2=|d\stackrel{}{z}|^2\frac{1}{6}|\stackrel{}{z}{}_{}{}^{}d\stackrel{}{z}|^2.$$
(5.4)
The corresponding vielbein is<sup>4</sup><sup>4</sup>4We inserted the $``$ sign in $`e^5`$ for later convenience.
$`e^1`$ $`=`$ $`\sqrt{{\displaystyle \frac{f_0}{3}}}\mathrm{cos}\theta \sigma ^1,`$ (5.5)
$`e^2`$ $`=`$ $`\sqrt{{\displaystyle \frac{f_0}{3}}}\sigma ^2,`$ (5.6)
$`e^3`$ $`=`$ $`{\displaystyle \frac{\sqrt{f_0}}{3}}\sqrt{3+\mathrm{cos}^2\theta }\left(\sigma ^3{\displaystyle \frac{4\mathrm{cos}\theta }{3+\mathrm{cos}^2\theta }}d\varphi \right),`$ (5.7)
$`e^4`$ $`=`$ $`\sqrt{{\displaystyle \frac{f_0}{3}}}d\theta ,`$ (5.8)
$`e^5`$ $`=`$ $`\sqrt{{\displaystyle \frac{f_0}{3}}}{\displaystyle \frac{2\mathrm{sin}\theta }{\sqrt{3+\mathrm{cos}^2\theta }}}d\varphi ,`$ (5.9)
with a warp factor
$$\mathrm{\Omega }^2=\sqrt{f_0}.$$
(5.10)
All the other fields are of course vanishing.
### 5.2 The Pilch-Warner fixed point solution
The vielbein in the Pilch-Warner solution is
$`e^1`$ $`=`$ $`\sqrt{{\displaystyle \frac{f_0}{3}}}\mathrm{cos}\theta \sigma ^1,`$ (5.11)
$`e^2`$ $`=`$ $`\sqrt{{\displaystyle \frac{f_0}{3}}}\mathrm{cos}\theta \sigma ^2,`$ (5.12)
$`e^3`$ $`=`$ $`{\displaystyle \frac{\sqrt{2f_0}}{3}}\mathrm{cos}\theta \sqrt{{\displaystyle \frac{3\mathrm{cos}^2\theta }{2\mathrm{cos}^2\theta }}}\left(\sigma ^3{\displaystyle \frac{2}{3\mathrm{cos}^2\theta }}d\varphi \right),`$ (5.13)
$`e^4`$ $`=`$ $`\sqrt{{\displaystyle \frac{2f_0}{3}}}\sqrt{2\mathrm{cos}^2\theta }d\theta ,`$ (5.14)
$`e^5`$ $`=`$ $`2\sqrt{{\displaystyle \frac{f_0}{3}}}\mathrm{sin}\theta \sqrt{{\displaystyle \frac{2\mathrm{cos}^2\theta }{3\mathrm{cos}^2\theta }}}d\varphi `$ (5.15)
and the warp factor is
$$\mathrm{\Omega }^2=\sqrt{f_0}\sqrt{2\mathrm{cos}^2\theta }.$$
(5.16)
Again one has complete coverage of the $`S^5`$ by the action of the $`SU(2)\times U(1)`$ if one takes $`0\theta \frac{\pi }{2}`$. The $`_2`$ that reduces the manifold to $`S^5/_2`$ lives inside the $`SU(2)`$ and so does not change the range of $`\theta `$.
This set of frames differs from the one in by a shift $`\sigma ^3\sigma ^32d\varphi `$. This shift is useful so as to make the assignment of the four-dimensional R-charge more transparent. In this frame the three-form flux is invariant under the four-dimensional R-symmetry.
### 5.3 Realization of the $`_2`$ symmetry
The theory of a single D3-brane probe is the reduction of the $`SU(N)\times SU(N)`$ gauge theory to a gauge theory with a single diagonal $`U(1)`$. One can see that the $`_2`$ symmetry acts as
$$(z_1,z_2,z_3,z_4)(z_1,z_2,z_3,z_4)$$
(5.17)
on the geometry. This corresponds to a shift $`\phi _3\phi _3+\pi `$ in the third Euler angle. This symmetry preserves the $`SO(3)`$ orbits. Based upon the field theory analysis, we expect the interpolating solutions to have the same property, i.e. the geometric action of the $`_2`$ symmetry will be implemented in the same way.
The action of the diffeomorphism on the vielbein is
$$(e^1,e^2,e^3,e^4,e^5)(e^1,e^2,e^3,e^4,e^5).$$
(5.18)
This can be undone by a local Lorentz rotation in the 1-2 plane by $`\pi `$. Since in the field theory the $`_2`$ symmetry is a charge conjugation, the type IIB realization has to contain $`S^2`$, which is world sheet orientation reversal. However, this acts on the $`SL(2,\mathrm{IR})/U(1)`$ coset fields as $`V_\alpha ^\pm V_\alpha ^\pm `$. This has to be undone by a type IIB R-symmetry rotation by $`\pi `$.
The Pilch-Warner solution respects the same $`_2`$ symmetry, which is consistent with the $`_2`$ action in the field theory dual. This further supports the expectation that this $`_2`$ will indeed be a symmetry of the complete interpolating family.
## 6 The interpolating solutions
### 6.1 Restrictions of the symmetries on the Ansatz
The most general five-dimensional metric respecting the $`SU(2)\times U(1)_R`$ symmetry is an $`SU(2)\times U(1)`$ fibration over an interval. Using coordinate reparametrization invariance in the fiber directions, this can be brought into the form
$`e^1`$ $`=`$ $`A_1(\sigma ^1+C_1d\varphi +C_2d\theta ),`$ (6.1)
$`e^2`$ $`=`$ $`A_2(\sigma ^2+D_1d\varphi +D_2d\theta ),`$ (6.2)
$`e^3`$ $`=`$ $`A_3(\sigma ^3+B_1d\varphi +B_2d\theta ),`$ (6.3)
$`e^4`$ $`=`$ $`A_4d\theta ,`$ (6.4)
$`e^5`$ $`=`$ $`A_5d\varphi ,`$ (6.5)
However, under the $`_2`$ symmetry $`\sigma ^1`$ and $`\sigma ^2`$ are odd, whereas $`\sigma ^3`$, $`d\theta `$ and $`d\varphi `$ are invariant. This constrains the Ansatz to
$`e^1`$ $`=`$ $`A_1\sigma ^1,`$ (6.6)
$`e^2`$ $`=`$ $`A_2\sigma ^2,`$ (6.7)
$`e^3`$ $`=`$ $`A_3(\sigma ^3+B_1d\varphi +B_2d\theta ),`$ (6.8)
$`e^4`$ $`=`$ $`A_4d\theta ,`$ (6.9)
$`e^5`$ $`=`$ $`A_5d\varphi .`$ (6.10)
In Appendix B we give the components of the spin connection for this metric.
The most general Ansatz for the three-form flux, $`G`$, that respects all the symmetries is:
$$\begin{array}{cc}\hfill G& =g_1(e^{134}ie^{234})+g_2(e^{145}ie^{245})+g_3(e^{135}ie^{235})+\hfill \\ & +g_4(e^{134}+ie^{234})+g_5(e^{145}+ie^{245})+g_6(e^{135}+ie^{235})\hfill \end{array}$$
(6.11)
and the most general dilaton-axion background respecting all the symmetries is
$$P=pe^4\mathrm{and}Q=0.$$
(6.12)
Note, that the $`U(1)`$ connection $`Q`$ has been gauged away.
The $`_2`$ symmetry acts through the diffeomorphism, the local Lorentz rotation by $`\pi `$ and a ten-dimensional R-symmetry rotation by $`\pi `$ on the Killing spinor. This imposes a projector on the Killing spinor
$$i\gamma ^{12}\chi ^{(1)}=\chi ^{(1)}\mathrm{and}i\gamma ^{12}\chi ^{(2)}{}_{}{}^{}=\chi ^{(2)}{}_{}{}^{}.$$
(6.13)
This projection restricts the spinors $`\chi ^{(1)}`$ and $`\chi ^{(2)}^{}`$ to live in the same two-dimensional subspace of the four-dimensional spinor space.
### 6.2 Solving the supersymmetry variations
#### 6.2.1 The “Magical Combination”
The magical combination $`2\delta \psi _0^{(\eta )}+\gamma ^1\delta \psi _1^{(\eta )}+\gamma ^2\delta \psi _2^{(\eta )}`$ of the gravitino variation equations is independent of all the fluxes, and depends only upon the metric. This leads to the projector equations:
$`\left({\displaystyle \frac{(A_1A_2)^{}}{A_4}}\gamma ^4iA_3\gamma ^3\right)\chi ^{(1)}`$ $`=`$ $`{\displaystyle \frac{2iA_1A_2}{\mathrm{\Omega }^2}}\chi ^{(1)},`$ (6.14)
$`({\displaystyle \frac{(A_1A_2)^{}}{A_4}}\gamma ^4iA_3\gamma ^3)\chi ^{(2)}^{}`$ $`=`$ $`{\displaystyle \frac{2iA_1A_2}{\mathrm{\Omega }^2}}\chi ^{(2)}^{}`$ (6.15)
In order for the foregoing projector equations to have non-trivial solutions, the metric coefficients must satisfy the condition:
$$A_3^2=\left(\frac{(A_1A_2)^{}}{A_4}\right)^2+\left(\frac{2A_1A_2}{\mathrm{\Omega }^2}\right)^2.$$
(6.16)
This condition is equivalent to setting:
$$\frac{(A_1A_2)^{}}{A_3A_4}=\mathrm{cos}\alpha \mathrm{and}\frac{2A_1A_2}{\mathrm{\Omega }^2A_3}=\mathrm{sin}\alpha ,$$
(6.17)
for some function, $`\alpha (\theta )`$. The Killing spinors then take the form
$`\chi ^{(1)}`$ $`=`$ $`\beta _1e^{\frac{i}{2}\varphi }(\mathrm{sin}{\displaystyle \frac{\alpha }{2}}|++\mathrm{cos}{\displaystyle \frac{\alpha }{2}}|+),`$ (6.18)
$`\chi ^{(2)}^{}`$ $`=`$ $`\beta _2^{}e^{\frac{i}{2}\varphi }(\mathrm{sin}{\displaystyle \frac{\alpha }{2}}|+++\mathrm{cos}{\displaystyle \frac{\alpha }{2}}|+),`$ (6.19)
where the $`\pm `$’s refer to the helicities of $`i\gamma ^{12}`$ and $`i\gamma ^{34}`$ on $`X_4`$. For consistency of the projector equation, the metric coefficients and the function $`\alpha `$ have to satisfy the differential equation
$$\frac{1}{2}(\mathrm{\Omega }^2A_3\mathrm{sin}\alpha )^{}=A_3A_4\mathrm{cos}\alpha .$$
(6.20)
We will assume in the following that for the interpolating solutions both spinors $`\chi ^{(1)}`$ and $`\chi ^{(2)}^{}`$ are non-vanishing and for this reason are linearly independent.
#### 6.2.2 The normalization conditions
After exploiting the second projector equation, we use normalization conditions for the Killing spinors coming from the symmetry algebra of the problem. The coefficients $`\beta _1`$ and $`\beta _2`$ have to satisfy the normalization condition (4.8):
$$|\beta _1|^2+|\beta _2|^2=1.$$
(6.21)
The other nornalization conditions (4.9) lead to the equations
$$\frac{3A_5}{2\mathrm{\Omega }^2}=\mathrm{cos}\alpha \mathrm{and}\frac{3A_3B_1}{2\mathrm{\Omega }^2}=\mathrm{sin}\alpha (|\beta _1|^2|\beta _2|^2).$$
(6.22)
For a range of $`0\alpha \frac{\pi }{2}`$ the vielbein coefficient $`A_5`$ has to be negative<sup>5</sup><sup>5</sup>5Note that this is just a convention and $`\alpha `$ can also be chosen in the range $`\frac{\pi }{2}\alpha \pi `$..
#### 6.2.3 The dilatino variation
The vanishing of the dilatino variations implies
$`g_4`$ $`=`$ $`{\displaystyle \frac{p}{\mathrm{\Omega }^2}}\left({\displaystyle \frac{\beta _2}{\beta _1}}{\displaystyle \frac{\beta _1^{}}{\beta _2^{}}}\right),`$ (6.23)
$`g_5`$ $`=`$ $`{\displaystyle \frac{p}{\mathrm{\Omega }^2\mathrm{tan}\alpha }}\left({\displaystyle \frac{\beta _2}{\beta _1}}+{\displaystyle \frac{\beta _1^{}}{\beta _2^{}}}\right),`$ (6.24)
$`ig_6`$ $`=`$ $`{\displaystyle \frac{p}{\mathrm{\Omega }^2\mathrm{sin}\alpha }}\left({\displaystyle \frac{\beta _2}{\beta _1}}+{\displaystyle \frac{\beta _1^{}}{\beta _2^{}}}\right).`$ (6.25)
Note that all three expressions have the same phase. This observation is important for the reality conditions.
#### 6.2.4 Reality conditions
The next big simplification of the problem comes from realizing that the fermion variation equations imply strong reality constraints. This is due to the reality of all the coefficients in the vielbein. The external gravitino variation equations imply
$$\mathrm{Im}\left(\frac{\beta _2}{\beta _1}g_1\right)=\mathrm{Im}\left(\frac{\beta _2}{\beta _1}g_2\right)=\mathrm{Re}\left(\frac{\beta _2}{\beta _1}g_3\right)=0,$$
(6.26)
and the “anti-magical” combination $`\gamma ^1\delta \psi _1^{(\eta )}\gamma ^2\delta \psi _2^{(\eta )}`$ of gravitino variations implies
$$B_2=0.$$
(6.27)
The gravitino variation equation $`\delta \psi _4^{(1)}`$ then turns into two differential equations for $`\beta _1`$, which take the form
$$a_{s,1}\frac{\beta _1^{}}{\beta _1}+a_{s,2}+a_{s,3}\frac{\beta _2}{\beta _1}=0,s=1,2,$$
(6.28)
with real coefficients $`a_{s,t}`$. One can take a linear combination of those two equations such that the term proportional to $`\frac{\beta _2}{\beta _1}`$ vanishes. This implies that the phases of $`\beta _1`$ and $`\beta _2`$ do not depend on $`\theta `$.
One can use the ten-dimensional $`U(1)`$ R-symmetry to give the same phase to $`\beta _1`$ and $`\beta _2`$. In addition one can multiply the spinors $`\chi ^{(\eta )}`$ by an arbitrary constant phase. This allows one to take $`\beta _1`$ and $`\beta _2`$ to be real, and they can be written as:
$$\beta _1=\mathrm{cos}\frac{\beta }{2}\mathrm{and}\beta _2=\mathrm{sin}\frac{\beta }{2}.$$
(6.29)
With this form of $`\beta _1,\beta _2`$, the spinor Ansatz in (6.18) and (6.19) is equivalent to introducing a dielectric projector as in .
It also follows that $`g_1`$ and $`g_2`$ are real, $`g_3`$ is imaginary and from the anti-magical combination it follows that $`p`$ is real. This means that all the complex functions in the problem become real functions.
In order to proceed with the gravitino variation equations it is useful to define the matrices
$$\begin{array}{c}L=\mathrm{cot}\alpha \gamma ^4i\mathrm{csc}\alpha \gamma ^3,\hfill \\ M=\frac{1}{\mathrm{sin}\beta \mathrm{sin}\alpha }\gamma ^1(1\mathrm{l}+\mathrm{cos}\beta \mathrm{sin}\alpha \gamma ^3+i\mathrm{cos}\alpha \gamma ^{34}),\hfill \\ N=\left(\frac{\mathrm{cos}\beta \beta ^{}}{2\mathrm{sin}\beta }+\frac{\mathrm{cot}\alpha \alpha ^{}}{2}\right)1\mathrm{l}+\frac{\beta ^{}}{2\mathrm{sin}\beta \mathrm{sin}\alpha }\gamma ^3+\frac{i\mathrm{cot}\alpha \beta ^{}}{2\mathrm{sin}\beta }\gamma ^4+\frac{i\alpha ^{}}{2\mathrm{sin}\alpha }\gamma ^{34}.\hfill \end{array}$$
(6.30)
These matrices satisfy the identities:
$`i\chi ^{(1)}=L\chi ^{(1)}`$ $`\mathrm{and}`$ $`i\chi ^{(2)}{}_{}{}^{}=L\chi ^{(2)}{}_{}{}^{},`$ (6.31)
$`\chi ^{(2)}=M\chi ^{(1)}`$ $`\mathrm{and}`$ $`\chi ^{(1)}{}_{}{}^{}=M\chi ^{(2)}{}_{}{}^{},`$ (6.32)
$`\chi ^{(1)}{}_{}{}^{}=N\chi ^{(1)}`$ $`\mathrm{and}`$ $`\chi ^{(2)}{}_{}{}^{}{}_{}{}^{}=N\chi ^{(2)}{}_{}{}^{}.`$ (6.33)
This enables one to rewrite the gravitino variation equations in the form:
$$R\chi ^{(1)}=0\mathrm{and}R\chi ^{(2)}{}_{}{}^{}=0,$$
(6.34)
for some matrix, $`R`$. This implies that $`R=0`$ modulo $`i\gamma ^{12}=1\mathrm{l}`$, and so one can read off the gravitino variation equations as the coefficients of $`1\mathrm{l}`$, $`\gamma ^3`$, $`\gamma ^4`$, $`\gamma ^{34}`$.
#### 6.2.5 The external gravitino variation
The external gravitino variation equations can be solved for $`\mathrm{\Omega }`$, $`g_1`$, $`g_2`$ and $`g_3`$
$`{\displaystyle \frac{f_0}{\mathrm{\Omega }^4}}`$ $`=`$ $`\mathrm{cos}\beta ,`$ (6.35)
$`g_1`$ $`=`$ $`{\displaystyle \frac{4\mathrm{cot}\beta \mathrm{\Omega }^{}}{\mathrm{\Omega }^3A_4}}={\displaystyle \frac{\beta ^{}}{\mathrm{\Omega }^2A_4}},`$ (6.36)
$`g_2`$ $`=`$ $`{\displaystyle \frac{4\mathrm{sin}\beta }{\mathrm{\Omega }^4}}{\displaystyle \frac{4\mathrm{cot}\alpha \mathrm{\Omega }^{}}{\mathrm{\Omega }^3A_4\mathrm{sin}\beta }}={\displaystyle \frac{4\mathrm{sin}\beta }{\mathrm{\Omega }^4}}{\displaystyle \frac{\mathrm{cot}\alpha \beta ^{}}{\mathrm{\Omega }^2A_4\mathrm{cos}\beta }},`$ (6.37)
$`g_3`$ $`=`$ $`{\displaystyle \frac{4i\mathrm{\Omega }^{}}{\mathrm{sin}\beta \mathrm{sin}\alpha \mathrm{\Omega }^3A_4}}={\displaystyle \frac{i\beta ^{}}{\mathrm{cos}\beta \mathrm{sin}\alpha \mathrm{\Omega }^2A_4}}.`$ (6.38)
We will use these expressions for the three-form flux and $`\mathrm{\Omega }^4`$ to simplify the remaining gravitino variation equations.
#### 6.2.6 The “anti-magical combination”
The anti-magical combination $`\gamma ^1\delta \psi _1^{(\eta )}\gamma ^2\delta \psi _2^{(\eta )}`$ leads to the following equations
$`{\displaystyle \frac{A_1}{A_2A_3}}{\displaystyle \frac{A_2}{A_1A_3}}`$ $`=`$ $`{\displaystyle \frac{4p\mathrm{cos}\alpha }{\mathrm{sin}^2\beta \mathrm{sin}^2\alpha }},`$ (6.39)
$`{\displaystyle \frac{A_1^{}}{A_1A_4}}{\displaystyle \frac{A_2^{}}{A_2A_4}}`$ $`=`$ $`2p\left(1{\displaystyle \frac{2}{\mathrm{sin}^2\beta \mathrm{sin}^2\alpha }}\right),`$ (6.40)
$`{\displaystyle \frac{A_1B_1}{A_2A_5}}{\displaystyle \frac{A_2B_1}{A_1A_5}}`$ $`=`$ $`{\displaystyle \frac{4p\mathrm{cos}\beta }{\mathrm{sin}^2\beta \mathrm{sin}\alpha }}.`$ (6.41)
Using the normalization conditions, (6.22), one can see that the last equation is actually redundant. Those are the only gravitino variations that contain $`p`$, but no $`g_1`$, $`g_2`$ or $`g_3`$. All the other gravitino variations do not contain $`p`$. A vanishing $`p`$ would imply that $`A_1=A_2`$, which inevitably leads to the Pilch-Warner solution.
#### 6.2.7 The gravitino variations in the third direction
$`{\displaystyle \frac{A_1}{A_2A_3}}+{\displaystyle \frac{A_2}{A_1A_3}}{\displaystyle \frac{A_3}{A_1A_2}}`$ $`=`$ $`{\displaystyle \frac{2}{\mathrm{\Omega }^2\mathrm{sin}\alpha }}{\displaystyle \frac{2\mathrm{cos}\alpha \beta ^{}}{A_4\mathrm{sin}\beta \mathrm{cos}\beta \mathrm{sin}^2\alpha }},`$ (6.42)
$`{\displaystyle \frac{A_3B_1^{}}{2A_4A_5}}`$ $`=`$ $`{\displaystyle \frac{2\mathrm{cos}\beta }{\mathrm{\Omega }^2}}+{\displaystyle \frac{\mathrm{cot}\alpha \beta ^{}}{A_4\mathrm{sin}\beta }},`$ (6.43)
$`{\displaystyle \frac{A_3^{}}{A_3A_4}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{cot}\alpha }{\mathrm{\Omega }^2}}{\displaystyle \frac{(2\mathrm{cos}^2\alpha +\mathrm{sin}^2\beta \mathrm{sin}^2\alpha )\beta ^{}}{2A_4\mathrm{sin}\beta \mathrm{cos}\beta \mathrm{sin}^2\alpha }},`$ (6.44)
#### 6.2.8 The gravitino variations in the fourth direction
$`{\displaystyle \frac{\alpha ^{}}{A_4}}`$ $`=`$ $`{\displaystyle \frac{3}{\mathrm{\Omega }^2}}+{\displaystyle \frac{\mathrm{cot}\alpha \beta ^{}}{A_4\mathrm{sin}\beta \mathrm{cos}\beta }},`$ (6.45)
$`{\displaystyle \frac{A_3B_1^{}}{2A_4A_5}}`$ $`=`$ $`{\displaystyle \frac{2\mathrm{cos}\beta }{\mathrm{\Omega }^2}}+{\displaystyle \frac{\mathrm{cos}\alpha \beta ^{}}{A_4\mathrm{sin}\beta \mathrm{sin}\alpha }}.`$ (6.46)
#### 6.2.9 The gravitino variations in the fifth direction
$`{\displaystyle \frac{A_1B_1}{A_2A_5}}+{\displaystyle \frac{A_2B_1}{A_1A_5}}`$ $`=`$ $`{\displaystyle \frac{2\beta ^{}}{A_4\mathrm{sin}\beta \mathrm{sin}\alpha }},`$ (6.47)
$`{\displaystyle \frac{A_5^{}}{A_5A_4}}`$ $`=`$ $`{\displaystyle \frac{2}{A_5\mathrm{sin}\alpha }}+{\displaystyle \frac{3\mathrm{cot}\alpha }{\mathrm{\Omega }^2}}{\displaystyle \frac{(2\mathrm{sin}^2\beta )\beta ^{}}{2A_4\mathrm{sin}\beta \mathrm{cos}\beta }},`$ (6.48)
$`{\displaystyle \frac{2\mathrm{cos}\alpha }{A_5}}`$ $`=`$ $`{\displaystyle \frac{3}{\mathrm{\Omega }^2}},`$ (6.49)
$`{\displaystyle \frac{A_3B_1^{}}{2A_4A_5}}`$ $`=`$ $`{\displaystyle \frac{2\mathrm{cos}\beta }{\mathrm{\Omega }^2}}+{\displaystyle \frac{\mathrm{cot}\alpha \beta ^{}}{A_4\mathrm{sin}\beta }}.`$ (6.50)
The third equation is equivalent to one of the normalization conditions. This confirms that the normalization conditions (6.22) are chosen with the correct normalization constant.
### 6.3 The BPS equations
One can eliminate most variables from the BPS equations and the normalization conditions (6.22). This leaves three independent equations for $`\alpha `$, $`\beta `$, $`\frac{A_1}{A_2}`$ and $`\frac{A_4}{\mathrm{\Omega }^2}`$. For notational simplicity we define
$$g=\frac{A_4}{\mathrm{\Omega }^2}\mathrm{and}h=\frac{A_1}{A_2}.$$
(6.51)
With these definitions, the BPS equations are:
$`\left(\mathrm{log}\left({\displaystyle \frac{g\mathrm{sin}\beta \mathrm{sin}^3\alpha (h+h^1)}{\mathrm{cos}\alpha \beta ^{}}}\right)\right)^{}`$ $`=`$ $`2g\mathrm{cot}\alpha ,`$ (6.52)
$`\left(\mathrm{log}{\displaystyle \frac{hh^1}{\mathrm{cos}\beta }}\right)^{}`$ $`=`$ $`{\displaystyle \frac{2\beta ^{}}{\mathrm{sin}\beta \mathrm{cos}\beta \mathrm{sin}^2\alpha }},`$ (6.53)
$`\left(\mathrm{log}{\displaystyle \frac{\mathrm{cot}\beta }{\mathrm{cos}\alpha }}\right)^{}`$ $`=`$ $`3g\mathrm{tan}\alpha .`$ (6.54)
It is straightforward to verify that BPS equations imply the supersymmetries, the Bianchi identities and the equations of motion. Once one has a solution to this system one can obtain every other field from $`g,h`$ and $`\beta `$. In Appendix C we have summarized all the equations needed to achieve this.
One can write (6.52)–(6.54) as a strictly first-order system by solving (6.53) for $`\beta ^{}`$ and substituting the results into (6.52) to obtain:
$$\left(\mathrm{log}\left(\frac{g(h+h^1)}{(3g1)}\frac{\mathrm{sin}^2\alpha }{\mathrm{cos}\beta }\right)\right)^{}=2g\mathrm{cot}\alpha .$$
(6.55)
It is also convenient to use this to substitute for $`\beta ^{}`$ on the right-hand side of (6.53) to arrive at:
$$\left(\mathrm{log}\frac{hh^1}{\mathrm{sin}^2\alpha \mathrm{cos}\beta }\right)^{}=\frac{6g}{\mathrm{sin}\alpha \mathrm{cos}\alpha }.$$
(6.56)
We may then take the BPS system to be (6.54)–(6.56), and from this we see that there is now at least one obvious integral of motion that can be obtained by taking a simple linear combination of (6.54)–(6.56) so as to get zero on the right-hand side. Indeed,
$$_0\frac{g^3}{(3g1)^3}(hh^1)(h+h^1)^3\frac{\mathrm{sin}^4\alpha }{\mathrm{sin}^2\beta \mathrm{cos}^2\beta }$$
(6.57)
must be constant as a consequence of the BPS equations.
It is unclear whether this system of equations has a simple, closed form for its solution. The results from gauged supergravity suggest that there should be an explicit solution, but it has so far eluded us. In the next section we will discuss the two known (KW and PW) solutions and use numerical methods to show that the BPS equations lead to a family that interpolates between these two solutions.
## 7 Solving the BPS equations
We will not be able to find the general solution to the BPS equations. However, we establish the existence of a one parameter family of solutions in several different ways. For this purpose it is useful to first understand the boundary conditions. This will allow us to count the integration constants of the BPS equations. We will find the linear perturbation around the $`T^{1,1}`$ and Pilch-Warner fixed point solutions. Furthermore we find the solutions numerically.
### 7.1 Boundary conditions
The interpolating solutions are given by $`\mathrm{IRIP}^3\times S^1`$ fibrations over an interval. Since the family of solutions should involve trading flux for the Kähler modulus of the blow-up, the generic member of the family should have the same topology as $`T^{1,1}`$. The size of the two $`S^2`$’s will change as the three-form flux is changed, but the topology will only degenerate to the orbifold when one reaches the PW solution. This means that the generic member of the family of solutions should have exactly the same boundary conditions on the interval as the $`T^{1,1}`$ metric. That is, $`A_1`$ should vanish at one end of the $`\theta `$-interval and $`A_5`$ should vanish at the other end. This will then properly fix the topology of the $`\mathrm{IRIP}^3\times S^1`$ fibration. Note that the PW solution also satisfies these boundary conditions, and furthermore $`A_1,A_2`$ and $`A_3`$ all vanish at $`\theta =\frac{\pi }{2}`$. The vanishing of these extra metric functions merely reflects the collapsed two-cycle in the orbifold.
We also have not yet fixed the reparametrization invariance ($`\theta \stackrel{~}{\theta }(\theta )`$). We do this by requiring that $`\alpha `$, defined in (6.22), be the independent variable and we will adopt this choice henceforth. As we will show below, one has $`\alpha [0,\frac{\pi }{2}]`$.
Consider the end of the interval where $`A_5`$ vanishes and where, generically, the coefficients $`A_1,\mathrm{},A_4`$ and $`\mathrm{\Omega }^2`$ are finite. The coefficient $`B_1`$ is also generically non-vanishing as the Klebanov-Witten limit suggests. Then equation (6.22) implies that $`\alpha \frac{\pi }{2}`$. Assuming that $`g`$ is generic, equation (6.54) implies that $`\beta 0`$ and equation (6.56) implies that $`h1`$. Assuming that
$$\beta \left(\alpha \frac{\pi }{2}\right)^s\mathrm{and}hh^1\left(\alpha \frac{\pi }{2}\right)^t\mathrm{with}s,t>0$$
(7.1)
equation (6.54) implies $`s=3g1`$ and equation (6.53) implies $`t=2s`$. Equation (6.52) is then trivially satisfied in this limit.
The solution is regular if there is no conical singularity and that the fluxes behave in a regular way. The vanishing circle at this end of the interval is given by the vector field
$$_\varphi B_1\sigma _3.$$
(7.2)
There is no deficit angle if the metric coefficients satisfy
$$\underset{\alpha \frac{\pi }{2}}{lim}\frac{B_1A_4}{A_5^{}}.$$
(7.3)
The two known solutions impliy that $`\frac{B_1A_4}{A_5^{}}1`$, and other values of this would correspond to different families of solutions. One can readily check that
$$\frac{B_1A_4}{A_5^{}}s.$$
(7.4)
and so we must have $`s=1`$, $`t=2`$ and $`g\frac{2}{3}`$. Note that the vanishing circle is not an isometry of the geometry. The fluxes can behave like scalars, vectors or two-forms in the 4-5 plane. Regularity of the fluxes requires
$$p0,g_2g_50,g_1ig_30\mathrm{and}g_4+ig_60.$$
(7.5)
It is easy to see that all of those regularity conditions follow from the behaviour of $`\alpha `$, $`\beta `$, $`g`$ and $`h`$
$$\alpha \frac{\pi }{2},\beta c_1\left(\alpha \frac{\pi }{2}\right),g\frac{2}{3}\mathrm{and}h1+c_2\left(\alpha \frac{\pi }{2}\right)^2.$$
(7.6)
Since $`\beta =0`$ at this end of the interval, the Killing spinors are of “Becker type” and so supersymmetric D3-brane probes should feel no force and this locus should be a moduli space for such probes.
At the other end of the interval $`A_1`$ must vanish and the coefficients $`A_2,\mathrm{},A_5`$ and $`\mathrm{\Omega }^2`$ are generically non-vanishing, and so one must have $`h0`$. Equation (6.17) implies that $`\alpha 0`$. Equation (6.54) implies $`\beta ^{}0`$. Assuming that $`\beta `$ and $`g`$ stay at generic finite values, equation (6.54) implies
$$\beta ^{}\alpha \mathrm{and}h\alpha ^s\mathrm{with}s>0.$$
(7.7)
Equations (6.52) and (6.53) then imply $`s=1`$ and $`g\frac{1}{2}`$.
The vanishing cycle at this end of the interval is generated by $`\sigma _1`$. Absence of a conical singularity requires
$$\frac{A_1^{}}{A_4}\pm 1,$$
(7.8)
The $`T^{1,1}`$ solution actually has $`\frac{A_1^{}}{A_4}1`$. The condition for the flux to be regular is
$$p0,g_1g_40\mathrm{and}A_2(g_2g_5)iA_3(g_3+g_6).$$
(7.9)
It is easy to see that all of those regularity conditions follow from the behaviour of $`\alpha `$, $`\beta ^{}`$, $`g`$ and $`h`$.
$$\alpha 0,\beta ^{}c_3\alpha ,g\frac{1}{2}\mathrm{and}hc_4\alpha ,$$
(7.10)
where
$$c_3=\underset{\alpha 0}{lim}\frac{\mathrm{sin}\beta \mathrm{cos}\beta }{2}.$$
(7.11)
At this end of the interval $`\beta `$ is generically non-zero and supersymmetric D3-brane probes should have a non-trivial potential. However, if they puff up into D5-branes by the dielectric effect, such branes might settle into a supersymmetric configuration in this part of the geometry.
It is at this end of the interval that the $`\mathrm{IR}P^3`$ degenerates into an $`S^2`$ of finite size unless $`c_4=h^{}=\mathrm{}`$, which happens in the PW limit.
### 7.2 Integration constants
Using $`\alpha `$ as the independent variable, we see that (6.53)–(6.55) is a first order system for three functions, $`g,h`$ and $`\beta `$. There are thus, naively, three constants of integration, which may be thought of as the initial values of these functions at one end of the interval. However, we saw in the last subsection that regularity of the solution imposes some constraints on these initial conditions: We derived the behaviour of $`\beta `$, $`g`$ and $`h`$ on both ends of the interval in such a way that the the solution is regular, has the desired toplogy, and the BPS equations are satisfied to leading order. On each side of the interval this left two integation constants $`c_1`$, $`c_2`$ at $`\alpha =\frac{\pi }{2}`$ and $`c_3`$, $`c_4`$ at $`\alpha =0`$. The complete solution space of the set of BPS equations is thus three dimensional and regularity at each end of the interval selects a two-dimensional subspace at each end. Two two-dimensional subspaces in three dimensions generically intersect in a one-dimensional subspace, and so there will be a (real) one-dimensional family of solutions that are regular at both ends of the interval.
One can refine this argument using the integral of motion, (6.57). As we will show below, $`_0`$ is given by a simple combination of $`c_1`$ and $`c_2`$, and by a simple combination $`c_3`$ and $`c_4`$. Choosing a value of $`_0`$ reduces the general solution space to a two-dimensional space and the regular solutions starting at each end of the interval to two one-dimensional subspaces. These subspaces generically intersect at a point, and so given a value of $`_0`$ one should expect a single solution that is regular at both ends of the $`\alpha `$-interval. Thus one expects the family of solutions we seek to be swept out by varying $`_0`$. As we will show below, the explicit numerical solutions precisely bear out this picture.
One should recall that we did, in fact, expect a complex one-dimensional family of solutions. The reduction to a real one-dimensional space came about via some of the gauge choices and rotations we made earlier. The real one-dimensional solution space can be complexified by reintroducing a constant phase $`e^{i\phi }`$ to the three-form flux and a phase $`e^{2i\phi }`$ to the dilaton $`P`$. The other two complex moduli of the solution are the integration constant $`\tau _0`$ for the gradient equation for the dilaton-axion and the two-form flux through the $`S^2`$ at $`\alpha =\pi `$.
### 7.3 The Klebanov-Witten limit
For the $`T^{1,1}`$ solution, one can use equation (6.17) to determine the angle $`\alpha `$ in terms of $`\theta `$ and then eliminate $`\theta `$. This leads to
$$\beta =0,g=\frac{2}{3+\mathrm{cos}^2\alpha }\mathrm{and}h=\frac{\sqrt{3}\mathrm{sin}\alpha }{\sqrt{3+\mathrm{cos}^2\alpha }}.$$
(7.12)
This limit looks somewhat singular because $`\beta =0`$. However, the ratio $`\frac{\beta ^{}}{\mathrm{sin}\beta }`$ is not singular. It can be calculated using equation (6.53)
$$\frac{\beta ^{}}{\mathrm{sin}\beta }=\frac{3\mathrm{cos}^2\alpha }{3+\mathrm{cos}^2\alpha }\mathrm{tan}\alpha .$$
(7.13)
It is then easy to see that the equations (6.52) and (6.54) are satisfied. The integral of motion, (6.57), diverges and corresponds to the singular limit, $`_0=\mathrm{}`$.
In order to see that the Klebanov-Witten limit is a smooth limit, one can do some linearized analysis. Because $`\beta `$ is vanishing, one can expand the BPS equation in $`\delta \left(\frac{\beta ^{}}{\mathrm{sin}\beta }\right)`$, $`\delta g`$ and $`\delta h`$. In these variables the linearized BPS equations turn into a second order system together with a first order equation
$$\delta \beta ^{}=\frac{3\mathrm{cos}^2\alpha }{3+\mathrm{cos}^2\alpha }\mathrm{tan}\alpha \delta \beta .$$
(7.14)
The obvious solution to the second order system is the trivial one. The linearized perturbation is then given by
$$\delta \beta =\frac{\mathrm{cos}\alpha }{3+\mathrm{cos}^2\alpha }\delta c.$$
(7.15)
where $`\delta c`$ is a (small) constant of integration. This solution satisfies all the boundary conditions, especially $`\delta \beta ^{}(\alpha =0)=0`$ and $`\delta \beta \left(\alpha =\frac{\pi }{2}\right)=0`$.
It is easy to derive the perturbation of the fields from this. To linear order, the metric and the warp factor remain unchanged and the dilaton-axion is still zero, however the three-form flux is given by:
$`\delta g_1={\displaystyle \frac{\mathrm{sin}\alpha (3\mathrm{cos}^2\alpha )}{2f_0(3+\mathrm{cos}^2\alpha )}}\delta c,`$ $`\delta g_4={\displaystyle \frac{\mathrm{cos}^2\alpha \mathrm{sin}\alpha }{f_0(3+\mathrm{cos}^2\alpha )}}\delta c,`$ (7.16)
$`\delta g_2={\displaystyle \frac{\mathrm{cos}\alpha (5+\mathrm{cos}^2\alpha )}{2f_0(3+\mathrm{cos}^2\alpha )}}\delta c,`$ $`\delta g_5={\displaystyle \frac{\mathrm{cos}^3\alpha }{f_0(3+\mathrm{cos}^2\alpha )}}\delta c,`$ (7.17)
$`\delta g_3={\displaystyle \frac{i(3\mathrm{cos}^2\alpha )}{2f_0(3+\mathrm{cos}^2\alpha )}}\delta c,`$ $`\delta g_6={\displaystyle \frac{i\mathrm{cos}^2\alpha }{f_0(3+\mathrm{cos}^2\alpha )}}\delta c.`$ (7.18)
The non-vanishing $`\delta g_4`$, $`\delta g_5`$ and $`\delta g_6`$ imply that at the quadratic order the dilaton-axion becomes non-trivial. It is easy to check that this perturbarion satisfies all the boundary conditions.
Since the $`T^{1,1}`$ solution has no three-form flux and has a trivial dilaton-axion background, it is invariant under the phase rotation $`e^{i\phi }`$. For this reason the perturbation can be complexified by complexifying $`\delta c`$.
### 7.4 The Pilch-Warner limit
At the Pilch-Warner fixed point one can show:
$$\mathrm{cos}\beta =\frac{3\mathrm{cos}^2\alpha }{3+\mathrm{cos}^2\alpha },\mathrm{sin}\beta =\frac{2\sqrt{3}\mathrm{cos}\alpha }{3+\mathrm{cos}^2\alpha },g=\frac{2}{3\mathrm{cos}^2\alpha }\mathrm{and}h=1.$$
(7.19)
Again, this limit looks somewhat singular, but equation (6.53) defines the derivative of the logarithm of a vanishing quantity.
$$\left(\mathrm{log}\left(hh^1\right)\right)^{}=\frac{18+14\mathrm{cos}^4\alpha }{\mathrm{sin}\alpha \mathrm{cos}\alpha (3\mathrm{cos}^2\alpha )(3+\mathrm{cos}^2\alpha )}$$
(7.20)
It is easy to check that the other two BPS equations are satisfied. The integral of motion, (6.57), has the value, $`_0=0`$.
As for the Klebanov-Witten limit, one can do a linearized analysis around the Pilch-Warner point. The BPS equations can be expanded in terms of $`\delta \beta `$, $`\delta g`$ and $`\delta \left(\mathrm{log}h\right)^{}`$. Again this leads to a second order system together with a first order equation
$$\delta h^{}=\frac{18+14\mathrm{cos}^4\alpha }{\mathrm{sin}\alpha \mathrm{cos}\alpha (3\mathrm{cos}^2\alpha )(3+\mathrm{cos}^2\alpha )}\delta h.$$
(7.21)
Again, the third order system can be solved by the trivial solution. The linearized perturbation is then given by
$$\delta h=\frac{\mathrm{cos}^2\alpha (3\mathrm{cos}^2\alpha )^2}{\mathrm{sin}^4\alpha (3+\mathrm{cos}^2\alpha )}\delta c,$$
(7.22)
where $`\delta c`$ is a (small) integration constant. This perturbation vanishes at $`\alpha =\frac{\pi }{2}`$ and diverges at $`\alpha =0`$. The divergence is due to the fact that this perturbation generates a resolution of the $`A_1`$ singularity in the Pilch-Warner geometry. A very similar behavior occurs if one perturbatively expands the resolution of the $`A_1`$ singularity in the Eguchi-Hanson geometry. This is discussed in Appendix D.
The non-vanishing perturbations of the vielbein coefficients are given by
$`\delta A_1`$ $`=`$ $`A_1{\displaystyle \frac{\delta h}{\sqrt{2}}},`$ (7.23)
$`\delta A_2`$ $`=`$ $`A_1{\displaystyle \frac{\delta h}{\sqrt{2}}}.`$ (7.24)
This solution has a similar behavior as the blowup of an $`A_1`$ singularity, however, it is geometrically not the same because there are non-zero fluxes and curvatures. The sign of $`\delta A_1`$ suggests that this perturbation makes $`A_1`$ vanish at $`\theta =\frac{\pi }{2}`$ whereas $`A_2`$ and $`A_3`$ stay finite. Also, the perturbations $`\delta p`$, $`\delta g_4`$, $`\delta g_5`$ and $`\delta g_6`$ are non-vanishing, which shows that the interpolating solutions indeed have a non-trivial dilaton-axion.
Since the Pilch-Warner fixed point solution has a non-trivial three-form flux, it is not invariant under the phase rotation $`e^{i\phi }`$. For this reason the foregoing perturbation can be complexified by
$$\delta g_i=ig_i\delta \phi .$$
(7.25)
### 7.5 The round $`S^5/_2`$
Another very simple solution to the BPS equations is given by:
$$\beta =0,g=1\mathrm{and}h=1.$$
(7.26)
It is easy to check that this is actually the round $`S^5/_2`$. The regularity of the metric at $`\alpha =\frac{\pi }{2}`$ implies that $`\varphi `$ has a periodicity of $`3\pi `$. For this reason volume integrals have an extra factor of $`\frac{3}{2}`$. This is important for the central charge calculations in the next section.
### 7.6 Numerical solutions
We now set about obtaining numerical solutions to the system of equations (6.52)–(6.54), and we will indeed see that this system of equations leads to a family of solutions that interpolates between the Pilch-Warner and $`T^{1,1}`$ geometries. As in the previous section, we fix the freedom to reparametrize the $`\theta `$-coordinate, $`\theta \stackrel{~}{\theta }(\theta )`$, by taking $`\alpha [0,\frac{\pi }{2}]`$ to be the independent variable. The next step is to use (6.53)–(6.55) to obtain expressions for $`g^{}`$, $`h^{}`$ and $`\beta ^{}`$ in terms $`g,h`$ and $`\beta `$. One can then employ a simple Euler method to get the numerical solution once one has specified “initial velocities” for $`g,h`$ and $`\beta `$. A priori there are three constants of integration, but as we described earlier, regularity reduces this to a one parameter family of solutions parametrized by the value of $`_0`$.
We find the solutions for the functions $`g,h`$ and $`\beta `$ by “shooting,” that is, we vary initial data at $`\alpha =0`$ and adjust it so as to hit the proper values at $`\alpha =\frac{\pi }{2}`$. In particular, we make use of the asymptotics given in (7.10) and (7.6). At $`\alpha =0`$ one has $`\beta ^{}=0`$ and $`h=0`$ and so the equations in (6.52) appear to be somewhat singular, however a careful series expansion about $`\alpha =0`$ leads to a regular expansion of all the undetermined functions, and one finds:
$`h`$ $`=`$ $`c_4\alpha +O(\alpha ^3),g={\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{16}}(4c_4^21)\alpha ^2+O(\alpha ^4),`$ (7.27)
$`\beta `$ $`=`$ $`\beta _0{\displaystyle \frac{1}{4}}\mathrm{sin}\beta _0\mathrm{cos}\beta _0\alpha ^2+O(\alpha ^4).`$ (7.28)
Similarly, at $`\alpha =\frac{\pi }{2}`$ one finds:
$`h`$ $`=`$ $`1+c_2\left(\alpha {\displaystyle \frac{\pi }{2}}\right)^2+O\left(\left(\alpha {\displaystyle \frac{\pi }{2}}\right)^4\right),`$ (7.29)
$`g`$ $`=`$ $`{\displaystyle \frac{2}{3}}+{\displaystyle \frac{1}{9}}(3c_1^22)\left(\alpha {\displaystyle \frac{\pi }{2}}\right)^2+O\left(\left(\alpha {\displaystyle \frac{\pi }{2}}\right)^4\right),`$ (7.30)
$`\beta `$ $`=`$ $`c_1\left(\alpha {\displaystyle \frac{\pi }{2}}\right){\displaystyle \frac{1}{6}}c_1(3c_1^2)\left(\alpha {\displaystyle \frac{\pi }{2}}\right)^3+O\left(\left(\alpha {\displaystyle \frac{\pi }{2}}\right)^5\right).`$ (7.31)
There are thus two free parameters at either end of the interval: $`\beta _0`$ and $`c_4`$ at $`\alpha =0`$ and $`c_1`$ and $`c_2`$ at $`\alpha =\frac{\pi }{2}`$. One can use the series expansions to check that the constant of the motion, (6.57), is given by
$$_0=\frac{1}{c_4^4\mathrm{sin}^2\beta _0\mathrm{cos}^2\beta _0}=\frac{128}{27}\frac{c_2}{c_1^2}.$$
(7.32)
It is simplest to shoot from $`\alpha =0`$ where the value of $`\beta _0`$ is chosen so as to select the particular member of the family of solutions and then the value of $`c_4`$ is adjusted so that one arrives at $`g=\frac{2}{3}`$ as $`\alpha \frac{\pi }{2}`$. We use the series expansion at $`\alpha =0`$ (evolved to fairly high order) to start the numerical solution, and then simply use an Euler method to generate the complete solution. By choosing $`c_4`$ to arrive at $`g=\frac{2}{3}`$ one finds that the asymptotic behavior of all the three functions obeys the proper asymptotics at $`\alpha =\frac{\pi }{2}`$.
The functions $`g,h`$ and $`\beta `$ for the KW and PW solutions are given by (7.12) and (7.19). In particular, observe that for $`\alpha =0`$ one has $`\beta _0=\frac{\pi }{3}`$ for the PW solution. We therefore found the numerical solutions for several values of $`\beta _0`$ in the range $`0\beta _0\frac{\pi }{3}`$. The results for $`\beta `$, $`g`$ and $`h`$ are plotted in Figures 1, 2 and 3. We have also plotted the exact results for the KW and PW solutions. The value of the integral of motion, (7.32), monotonically increases across the family from $`0`$ for the PW solution to infinity for the KW solution. It is clear from these graphs that the solutions to the BPS equations do indeed interpolate between the KW and PW solutions, and that there is a smooth family of solutions in which the flux of the PW solution is traded for a blowing-up of the non-trivial two-cycle.
We have focussed on the set of regular solutions to the BPS equations. As we noted earlier, there is a three-parameter family of solutions in general. Our numerical solutions show that the other solutions to the BPS equations can be characterized as solutions starting from $`g=\frac{1}{2}`$ at $`\alpha =0`$ and arriving at some arbitrary value of $`g`$ at $`\alpha =\frac{\pi }{2}`$. The absence of conical singularities required that $`(3g1)`$, and the Klebanov-Witten solution imposed $`(3g1)=1`$. However, there might be other interesting solutions to our BPS equations that are regular geometries but with different asymptotic values of $`g`$.
## 8 The central charge
As a final check on our results, we calculate the central charge of the family of solutions. It turns out that this is actually an exact calculation even though the exact solutions are not known. The central charge of the holographic dual gauge theory is proportional to the effective five-dimensional Newton constant .
### 8.1 Calculating the effective five-dimensional Newton constant
The effective five-dimensional Newton constant $`G_5`$ is given by
$$G_5=G_{10}\underset{X_5}{}\frac{1}{\mathrm{\Omega }^2}e^1e^2e^3e^4e^5$$
(8.1)
Using the vielbein Ansatz this can be reduced to
$$G_5=2(2\pi )^3G_{10}\underset{I}{}\frac{A_1A_2A_3A_4A_5}{\mathrm{\Omega }^2}𝑑\theta $$
(8.2)
Using the equations (6.17), (6.22) one can show that
$$G_5=\frac{2(2\pi )^3G_{10}}{3}\underset{I}{}(A_1^2A_2^2)^{}𝑑\theta .$$
(8.3)
For the family of $`𝒩=1`$ theories the boundary conditions imply that
$$G_5^{(IR)}=\frac{2(2\pi )^3G_{10}f_0^2}{27},$$
(8.4)
whereas for the $`𝒩=2`$ theory
$$G_5^{(UV)}=\frac{3(2\pi )^3G_{10}f_0^2}{48}.$$
(8.5)
Note that the factor 3 in the enumerator is due to the different periodicity of $`\varphi `$ as discussed in section 7.5. These formulas depend on the integration constant, $`f_0`$. This constant was introduced as the coefficient of the five-form flux. For this reason it is related to the number $`N_{D3}`$ of D3-branes, which is the rank of the gauge group. However, (8.4) already shows that the central charge of the dual field theory is constant across the entire family of solutions, and independent of the choice of the initial data for the BPS equations. To conclude that this implies that the central charge of the dual field theory is constant across the family one really needs to show that the parameter, $`f_0`$, represents the number of $`D3`$-branes present in the family of solutions. While this seems highly plausible, we will now prove it.
### 8.2 Calculating the rank of the gauge group
The Bianchi identity
$$dF^{(5)}=\frac{i}{8}GG^{}$$
(8.6)
implies that the five-form flux is not only sourced by D3-branes, but also by three-form flux. For this reason the total five-form flux cannot be used to determine the rank of the gauge group. The effect of the three-form flux can be subtracted as follows
$$N_{D3}=\underset{X_5}{}\left(F^{(5)}\frac{i}{16}ϵ_{\alpha \beta }A^\alpha F^\beta \right).$$
(8.7)
The five-form under the integral is not gauge invariant by itself, but the integral is gauge invariant.
To determine this integral, we need to relate the quantities appearing here to metric and field coefficients. The internal part of the field strength $`F^{(5)}`$ is given by
$$F_{int}^{(5)}=\frac{f_0A_1A_2A_3A_4A_5}{\mathrm{\Omega }^{10}}\sigma ^1\sigma ^2\sigma ^3d\theta d\varphi $$
(8.8)
The three-form flux $`F^{(3)}=F^1=(F^2)^{}`$ is related to $`G`$ by
$$G=ϵ_{\alpha \beta }V_+^\alpha F^\beta .$$
(8.9)
Using the identity
$$|V_+^2|^2|V_+^1|^2=1,$$
(8.10)
the foregoing relation can be inverted to yield
$$F^{(3)}=V_+^2{}_{}{}^{}G+V_+^1G^{}.$$
(8.11)
In our geometry $`G`$ has the form
$$\begin{array}{cc}\hfill G=& h_1\sigma ^1\sigma ^3d\theta +h_2\sigma ^2\sigma ^3d\theta +h_3\sigma ^1\sigma ^3d\varphi +\hfill \\ \hfill +& h_4\sigma ^2\sigma ^3d\varphi +h_5\sigma ^1d\theta d\varphi +h_6\sigma ^2d\theta d\varphi ,\hfill \end{array}$$
(8.12)
which implies, that $`F^{(3)}`$ has the form
$$\begin{array}{cc}\hfill F^{(3)}& =f_1\sigma ^1\sigma ^3d\theta +f_2\sigma ^2\sigma ^3d\theta +f_3\sigma ^1\sigma ^3d\varphi +\hfill \\ \hfill +& f_4\sigma ^2\sigma ^3d\varphi +f_5\sigma ^1d\theta d\varphi +f_6\sigma ^2d\theta d\varphi .\hfill \end{array}$$
(8.13)
The field strength, $`F^{(3)}`$, satisfies the Bianchi identity $`dF^{(3)}=0`$, which implies
$$f_5=f_4^{}\mathrm{and}f_6=f_3^{}.$$
(8.14)
A two-form potential $`A^{(2)}`$ for such an $`F^{(3)}`$ is then
$$A^{(2)}=f_1\sigma ^2d\theta +f_2\sigma ^1d\theta f_3\sigma ^2d\varphi +f_4\sigma ^1d\varphi .$$
(8.15)
This can be used to determine
$$ϵ_{\alpha \beta }A^\alpha F^\beta =2(f_1f_3^{}f_1^{}f_3+f_2f_4^{}f_2^{}f_4)\sigma ^1\sigma ^2\sigma ^3d\theta d\varphi $$
(8.16)
or
$$ϵ_{\alpha \beta }A^\alpha F^\beta =2(h_1h_3^{}h_1^{}h_3+h_2h_4^{}h_2^{}h_4)\sigma ^1\sigma ^2\sigma ^3d\theta d\varphi $$
(8.17)
This can be reexpressed in terms of the vielbein coefficients
$$\begin{array}{c}ϵ_{\alpha \beta }A^\alpha F^\beta =\\ =4A_3^2A_4A_5((g_1+g_4)(g_3+g_6)A_1^2+(g_1g_4)(g_3g_6)A_2^2)\sigma ^1\sigma ^2\sigma ^3d\theta d\varphi \end{array}$$
(8.18)
One can check that
$$\begin{array}{c}\frac{(A_1^2A_2^2\mathrm{cos}^2\beta )^{}}{3f_0}=\hfill \\ \frac{f_0A_1A_2A_3A_4A_5}{\mathrm{\Omega }^{10}}+\frac{iA_3^2A_4A_5((g_1+g_4)(g_3+g_6)A_1^2+(g_1g_4)(g_3g_6)A_2^2)}{144}.\hfill \end{array}$$
(8.19)
which implies that
$$N_{D3}=\frac{2(2\pi )^3}{3f_0}\underset{I}{}(A_1^2A_2^2\mathrm{cos}^2\beta )^{}𝑑\theta .$$
(8.20)
For the family of $`𝒩=1`$ theories this yields:
$$N_{D3}^{(IR)}=\frac{2(2\pi )^3f_0}{27}$$
(8.21)
and for the $`𝒩=2`$ theory this is
$$N_{D3}^{(UV)}=\frac{3(2\pi )^3f_0}{48}.$$
(8.22)
This enables us to express the effective five-dimensional Newton constant in terms of the rank of the gauge group
$$G_5^{(IR)}=\frac{27G_{10}N_{D3}^{(IR)}^2}{2(2\pi )^3}\mathrm{and}G_5^{(UV)}=\frac{48G_{10}N_{D3}^{(IR)}^2}{3(2\pi )^3}.$$
(8.23)
The ratio of the effective five-dimensional Newton constants is exactly the ratio of the central charges of the UV and the IR gauge theories
$$\frac{G_5^{(IR)}}{G_5^{(UV)}}=\frac{27}{32}=\frac{c_{(IR)}}{c_{(UV)}}.$$
(8.24)
Thus the family of solutions has precisely the correct central charge to be the duals of the family of fixed points predicted in .
## 9 Conclusions
We have found the long-sought family of $`AdS_5`$ vacuum solutions that interpolate between the $`T^{1,1}`$ compactification and the flux compactification of Pilch and Warner . This family of solutions is holographically dual to the family of $`𝒩=1^{}`$ IR fixed points that can be obtained by flowing from an $`𝒩=2`$, $`_2`$ quiver gauge theory. In the field theory, this family is parametrized by the ratio, $`m_1/m_2`$, of masses given to the chiral multiplets on each node of the quiver. In supergravity the difference of the masses, $`m_1m_2`$, is dual to the Kähler modulus of a non-trivial $`S^2`$, while the sum of the masses, $`m_1+m_2`$, is dual to a non-trivial, three-form field strength. Thus the family represents a kind of continuous geometric transition in which a Kähler deformation is traded for flux.
One of the surprises, and perhaps one of the reasons why this solution was not discovered earlier, is that the generic solution has a non-trivial dilaton. It is surprising because the dilaton background is trivial for the two previously know (KW and PW) solutions. There are obvious questions about whether there is any interesting physics to be learned from the non-trivial dilaton profiles. On the more mathematical side, it raises questions about the underlying geometric structure of these solutions. One of the important insights of was that the geometry of the PW solution, and indeed the flows to and around it , possessed an integrable complex structure, and indeed were “almost Calabi-Yau.” The non-trivial dilaton profile, and indeed the fact that it is real, seems to be at odds with the integrability of the complex structure. We have tried the obvious generalizations of the integrable complex structure found in and they fail to work here, and this failure perhaps explains the incompatibility of the complex structures, noted in , of the PW flow and of the Calabi-Yau metric that must underlie the KW flow. There is thus an interesting issue as to how to characterize the geometry of the interpolating family obtained here.
The system of BPS equations that we obtained were surprisingly complicated, also probably as a consequence of the non-trivial dilaton profile. This is all the more surprising in the light of the results of that led to the conjectured existence of the family of solutions. It was shown in that, from the perspective of five-dimensional, $`𝒩=4`$ gauged supergravity, all the vacuum solutions in the family, and indeed all the flows to them, were governed by exactly the same set of equations. The complete family, in five-dimensional supergravity, is swept out by the action of an $`SU(2)`$ symmetry. One would therefore, naively, expect an equally simple formulation in ten-dimensions. However, as was pointed out in , and as we see explicitly here, this sweeping out of the family involves some extremely non-trivial trading of very different geometric quantities in ten dimensions. It is certainly not the first time that a trivial symmetry in lower dimensions has led to subtle or profound effects in higher dimensions, and indeed the parallels between the present example and mirror symmetry are rather intriguing. It would certainly be very interesting to find how the symmetry that sweeps out the family acts in ten dimensions. This might be similar to the $`SL(2,\mathrm{IR})`$ action in . For this reason there should be a simpler form of our BPS equations and a way to solve them analytically. However, in string theory such a continous symmetry group of the supergravity will be broken down to a discrete duality group by solitonic excitations .
There is also the issue of the flow solutions: We have found the fixed points, but it would be very useful to find the family of flows from the quiver gauge theories to these fixed points. Finding these might also shed light upon the underlying geometric structure.
As a final comment, we found the family of solutions by a very careful analysis of the symmetries of the field theory. In particular, the discrete $`_2`$ symmetry in combination with the $`SO(3)`$ symmetry played a very significant role in fixing the metric Ansatz and in determining one of the supersymmetry projectors. We suspect that such a careful treatment of such discrete symmetries of will also give new insights into how to solve other open problems in holographic descriptions of field theories, especially for field theories related to $`𝒩=4`$ SYM.
Acknowledgments
This work is supported in part by funds provided by the DOE under grant number DE-FG03-84ER-40168. The work of NH is supported in part by a Fletcher Jones Graduate Fellowship from USC. Research at the Perimeter Institute is supported in part by funds from NSERC of Canada.
NH would like to thank the theory group at Stony Brook and ANU for hospitality. CR and NW would like to thank the Aspen Center for Physics in which part of the work was done.
We would like to thank Andy Brandhuber, Alex Buchel, Rich Corrado, Jerome Gauntlett, Jaume Gomis, Peter Mayr, Andrei Starinets and Nemani Suryanarayana for useful discussions.
## Appendix A Some Clifford algebra
### A.1 Generalities
The Clifford algebra is defined by the anticommutation relations
$$\{\gamma ^m,\gamma ^n\}=2\eta ^{mn},$$
(A.1)
where $`\eta ^{mn}=\eta ^m\delta ^{mn}`$. We choose a representation in which $`\sqrt{\eta ^m}\gamma ^m`$ is Hermitean<sup>6</sup><sup>6</sup>6By the square root we mean $`\sqrt{1}=1`$ and $`\sqrt{1}=i`$.. Given a complex structure, one can define the raising and lowering operators
$$\mathrm{\Gamma }^m=\sqrt{\eta ^{2m1}}\gamma ^{2m1}+i\sqrt{\eta ^{2m}}\gamma ^{2m},\mathrm{and}(\mathrm{\Gamma }^m)^{}=\sqrt{\eta ^{2m1}}\gamma ^{2m1}i\sqrt{\eta ^{2m}}\gamma ^{2m}.$$
(A.2)
Then the raising and lowering operators satisfy the following anticommutation relations:
$$\{\mathrm{\Gamma }^m,\mathrm{\Gamma }^n\}=\{(\mathrm{\Gamma }^m)^{},(\mathrm{\Gamma }^n)^{}\}=0\mathrm{and}\{\mathrm{\Gamma }^m,(\mathrm{\Gamma }^n)^{}\}=4\delta ^{mn}.$$
(A.3)
One can then define the fermion number operators
$$F^m=i\sqrt{\eta ^{2m1}}\gamma ^{2m1}\sqrt{\eta ^{2m}}\gamma ^{2m}=1\frac{1}{2}\mathrm{\Gamma }^m(\mathrm{\Gamma }^m)^{}=1+\frac{1}{2}(\mathrm{\Gamma }^m)^{}\mathrm{\Gamma }^m.$$
(A.4)
The chirality operator is then the product of all the Fermion number operators $`\gamma =F^1\mathrm{}F^n`$.
The Fermion number operators have eigenvalues $`\pm 1`$. The eigenvalues of the Fermion number operators can be used to label a basis of states. One can define a ground state $`|0`$ which is anihilated by all the lowering operators. It has Fermion number $`1`$ for all Fermion number operators. All other states can be gotten by applying raising operators. If one labels a state by $`|\nu _1,\mathrm{},\nu _n`$, then the raising and lowering operators act as follows:
$`|\nu _1,\mathrm{},+1,\mathrm{},\nu _n`$ $`=`$ $`{\displaystyle \frac{1}{2}}\nu _1\mathrm{}\nu _{m1}(\mathrm{\Gamma }^m)^{}|\nu _1,\mathrm{},1,\mathrm{},\nu _n,`$ (A.5)
$`|\nu _1,\mathrm{},1,\mathrm{},\nu _n`$ $`=`$ $`{\displaystyle \frac{1}{2}}\nu _1\mathrm{}\nu _{m1}\mathrm{\Gamma }^m|\nu _1,\mathrm{},+1,\mathrm{},\nu _n.`$ (A.6)
This defines the matrix elements of the gamma matrices. One can see that in this basis $`\mathrm{\Gamma }^m`$ is real. From this follows that
* The matrices $`\sqrt{\eta ^m}\gamma ^m`$ are Hermitean,
* The matrices $`\sqrt{\eta ^{2m1}}\gamma ^{2m1}`$ are symmetric and real and
* The matrices $`\sqrt{\eta ^{2m}}\gamma ^{2m}`$ are antisymmetric and imaginary.
In general there are matrices $`B`$, $`C`$ and $`D`$ such that
$`(\gamma ^m)^{}`$ $`=`$ $`\eta _BB\gamma ^mB^1,`$ (A.7)
$`(\gamma ^m)^{}`$ $`=`$ $`C\gamma ^mC^1,`$ (A.8)
$`(\gamma ^m)^t`$ $`=`$ $`\eta _BD\gamma ^mD^1,`$ (A.9)
where $`\eta _B=\pm 1`$ is a constant which is chosen (if possible) such that $`BB^{}=1`$. One can see that $`D=(B^{})^1C`$. Given a spinor $`ϵ`$, $`ϵ^{}=B^1ϵ^{}`$, $`\overline{ϵ}=ϵ^{}C`$ and $`\stackrel{~}{ϵ}=ϵ^tD`$ transform covariantly.
If $`BB^{}=1`$ one can impose the Majorana condition $`ϵ=B^1ϵ^{}`$. And if $`B`$ commutes with the chirality operator $`\gamma `$, one can impose the Majorana-Weyl condition.
In the following we collect useful Gamma matrix identities in various dimensions.
### A.2 $`Spin(1,9)`$
Chirality operator:
$$\gamma =\gamma ^{0\mathrm{}9}$$
(A.10)
Complex conjugation:
$$B=\gamma ^{013579}$$
(A.11)
$$B\gamma ^MB^1=(\gamma ^M)^{}$$
(A.12)
$$B\gamma _{(10)}B^1=\gamma _{(10)}^{}$$
(A.13)
$$BB^{}=1$$
(A.14)
Hermitean conjugation:
$$C=\gamma ^0$$
(A.15)
$$C\gamma ^MC^1=(\gamma ^M)^{}$$
(A.16)
Transpose:
$$D=(B^{})^1C=\gamma ^{13579}$$
(A.17)
$$D\gamma ^MD^1=(\gamma ^M)^t$$
(A.18)
### A.3 $`Spin(1,4)`$
Chirality operator:
$$\gamma ^4=\gamma ^{0123}$$
(A.19)
$$\gamma ^{01234}=1$$
(A.20)
Complex Conjugation:
$$B=\gamma ^{013}$$
(A.21)
$$B\gamma ^\mu B^1=(\gamma ^\mu )^{}$$
(A.22)
$$BB^{}=1$$
(A.23)
Hermitean conjugation:
$$C=\gamma ^0$$
(A.24)
$$C\gamma ^\mu C^1=(\gamma ^\mu )^{}$$
(A.25)
Transpose:
$$D=(B^{})^1C=\gamma ^{13}$$
(A.26)
$$D\gamma ^\mu D^1=(\gamma ^\mu )^t$$
(A.27)
### A.4 $`Spin(5)`$
Chirality operator:
$$\gamma ^5=\gamma ^{1234}$$
(A.28)
$$\gamma ^{12345}=1$$
(A.29)
Complex Conjugation:
$$B=\gamma ^{24}$$
(A.30)
$$B\gamma ^mB^1=(\gamma ^m)^{}$$
(A.31)
$$BB^{}=1$$
(A.32)
Hermitean conjugation:
$$C=1$$
(A.33)
$$C\gamma ^mC^1=(\gamma ^m)^{}$$
(A.34)
Transpose:
$$D=(B^{})^1C=\gamma ^{24}$$
(A.35)
$$D\gamma ^mD^1=(\gamma ^m)^t$$
(A.36)
It is easy to check that
$$\begin{array}{cc}B^1|++^{}=|,& B^1|+^{}=|+,\\ B^1|+^{}=|+,& B^1|^{}=|++.\end{array}$$
(A.37)
### A.5 Decomposition of a ten-dimensional spinor
We want to decompose spinors in ten-dimensional Minkowski space of mostly minus signature into four-dimensional and six-dimensional spinors. The gamma matrices can be decomposed as
$$\gamma _{(10)}^\mu =\left(\begin{array}{cc}0& \gamma _{(e)}^\mu 1_{(i)}\\ \gamma _{(e)}^\mu 1_{(i)}& 0\end{array}\right)\mathrm{and}\gamma _{(10)}^m=\left(\begin{array}{cc}0& 1_{(e)}\gamma _{(i)}^m\\ 1_{(e)}\gamma _{(i)}^m& 0\end{array}\right).$$
(A.38)
Note that the internal gamma matrices $`\gamma _{(i)}^m`$ have a $`+`$-signaturte.
The ten-dimensional chirality operator is given by
$$\gamma _{(10)}=\gamma _{(10)}^{0\mathrm{}9}=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),$$
(A.39)
the complex conjugation is given by
$$B_{(10)}=\left(\begin{array}{cc}B_{(e)}B_{(i)}& 0\\ 0& B_{(e)}B_{(i)}\end{array}\right)$$
(A.40)
and the hermitean conjugation is given by
$$C_{(10)}=\left(\begin{array}{cc}0& C_{(e)}C_{(i)}\\ C_{(e)}C_{(i)}& 0\end{array}\right).$$
(A.41)
## Appendix B The spin connection of the internal metric
The derivatives of the vielbein are
$`de^1`$ $`=`$ $`{\displaystyle \frac{A_1^{}}{A_1A_4}}e^4e^1+{\displaystyle \frac{A_1}{A_2A_3}}e^2e^3{\displaystyle \frac{A_1B_1}{A_2A_5}}e^2e^5{\displaystyle \frac{A_1B_2}{A_2A_4}}e^2e^4,`$ (B.1)
$`de^2`$ $`=`$ $`{\displaystyle \frac{A_2^{}}{A_2A_4}}e^4e^2{\displaystyle \frac{A_2}{A_1A_3}}e^1e^3+{\displaystyle \frac{A_2B_1}{A_1A_5}}e^1e^5+{\displaystyle \frac{A_2B_2}{A_1A_4}}e^1e^4,`$ (B.2)
$`de^3`$ $`=`$ $`{\displaystyle \frac{A_3^{}}{A_3A_4}}e^4e^3+{\displaystyle \frac{A_3}{A_1A_2}}e^1e^2+{\displaystyle \frac{A_3B_1^{}}{A_4A_5}}e^4e^5,`$ (B.3)
$`de^4`$ $`=`$ $`0,`$ (B.4)
$`de^5`$ $`=`$ $`{\displaystyle \frac{A_5^{}}{A_5A_4}}e^4e^5.`$ (B.5)
This leads to the following spin connection:
$`\omega _{114}`$ $`=`$ $`{\displaystyle \frac{A_1^{}}{A_1A_4}},`$ (B.6)
$`\omega _{123}`$ $`=`$ $`{\displaystyle \frac{A_1}{2A_2A_3}}+{\displaystyle \frac{A_2}{2A_1A_3}}+{\displaystyle \frac{A_3}{2A_1A_2}},`$ (B.7)
$`\omega _{124}`$ $`=`$ $`{\displaystyle \frac{A_1B_2}{2A_2A_4}}{\displaystyle \frac{A_2B_2}{2A_1A_4}},`$ (B.8)
$`\omega _{125}`$ $`=`$ $`{\displaystyle \frac{A_1B_1}{2A_2A_5}}{\displaystyle \frac{A_2B_1}{2A_1A_5}},`$ (B.9)
$`\omega _{224}`$ $`=`$ $`{\displaystyle \frac{A_2^{}}{A_2A_4}},`$ (B.10)
$`\omega _{213}`$ $`=`$ $`{\displaystyle \frac{A_1}{2A_2A_3}}+{\displaystyle \frac{A_2}{2A_1A_3}}{\displaystyle \frac{A_3}{2A_1A_2}},`$ (B.11)
$`\omega _{214}`$ $`=`$ $`{\displaystyle \frac{A_1B_2}{2A_2A_4}}{\displaystyle \frac{A_2B_2}{2A_1A_4}},`$ (B.12)
$`\omega _{215}`$ $`=`$ $`{\displaystyle \frac{A_1B_1}{2A_2A_5}}{\displaystyle \frac{A_2B_1}{2A_1A_5}},`$ (B.13)
$`\omega _{312}`$ $`=`$ $`{\displaystyle \frac{A_1}{2A_2A_3}}+{\displaystyle \frac{A_2}{2A_1A_3}}{\displaystyle \frac{A_3}{2A_1A_2}},`$ (B.14)
$`\omega _{334}`$ $`=`$ $`{\displaystyle \frac{A_3^{}}{A_3A_4}},`$ (B.15)
$`\omega _{345}`$ $`=`$ $`{\displaystyle \frac{A_3B_1^{}}{2A_4A_5}},`$ (B.16)
$`\omega _{412}`$ $`=`$ $`{\displaystyle \frac{A_1B_2}{2A_2A_4}}{\displaystyle \frac{A_2B_2}{2A_1A_4}},`$ (B.17)
$`\omega _{435}`$ $`=`$ $`{\displaystyle \frac{A_3B_1^{}}{2A_4A_5}},`$ (B.18)
$`\omega _{512}`$ $`=`$ $`{\displaystyle \frac{A_1B_1}{2A_2A_5}}{\displaystyle \frac{A_2B_1}{2A_1A_5}},`$ (B.19)
$`\omega _{534}`$ $`=`$ $`{\displaystyle \frac{A_3B_1^{}}{2A_4A_5}},`$ (B.20)
$`\omega _{554}`$ $`=`$ $`{\displaystyle \frac{A_5^{}}{A_5A_4}}.`$ (B.21)
## Appendix C Recovering the fields
Going through all the independent BPS equations one can recover the vielbein coefficients from $`\alpha `$, $`\beta `$, $`g`$ and $`h`$
$`\mathrm{\Omega }^2`$ $`=`$ $`\sqrt{{\displaystyle \frac{f_0}{\mathrm{cos}\beta }}},`$ (C.1)
$`A_1`$ $`=`$ $`\sqrt{{\displaystyle \frac{f_0\mathrm{sin}^3\alpha \mathrm{sin}\beta g(h^2+1)}{4\beta ^{}\mathrm{cos}\alpha }}},`$ (C.2)
$`A_2`$ $`=`$ $`\sqrt{{\displaystyle \frac{f_0\mathrm{sin}^3\alpha \mathrm{sin}\beta g(h^2+1)}{4\beta ^{}\mathrm{cos}\alpha h^2}}},`$ (C.3)
$`A_3`$ $`=`$ $`{\displaystyle \frac{\sqrt{f_0\mathrm{cos}\beta }\mathrm{sin}^2\alpha \mathrm{sin}\beta g(h^2+1)}{2\beta ^{}\mathrm{cos}\alpha h}},`$ (C.4)
$`A_4`$ $`=`$ $`g\sqrt{{\displaystyle \frac{f_0}{\mathrm{cos}\beta }}},`$ (C.5)
$`A_5`$ $`=`$ $`{\displaystyle \frac{2}{3}}\sqrt{{\displaystyle \frac{f_0}{\mathrm{cos}\beta }}}\mathrm{cos}\alpha ,`$ (C.6)
$`B_1`$ $`=`$ $`{\displaystyle \frac{4\beta ^{}\mathrm{cos}\alpha h}{3g\mathrm{sin}\alpha \mathrm{sin}\beta (h^2+1)}},`$ (C.7)
$`p`$ $`=`$ $`{\displaystyle \frac{\beta ^{}\mathrm{sin}\beta (h^21)}{2\sqrt{f_0\mathrm{cos}\beta }g(h^2+1)}},`$ (C.8)
$`g_1`$ $`=`$ $`{\displaystyle \frac{\mathrm{cos}\beta \beta ^{}}{f_0g}},`$ (C.9)
$`g_2`$ $`=`$ $`{\displaystyle \frac{4\mathrm{sin}\beta \mathrm{cos}\beta }{f_0}}{\displaystyle \frac{\mathrm{cot}\alpha \beta ^{}}{f_0g}},`$ (C.10)
$`g_3`$ $`=`$ $`{\displaystyle \frac{i\beta ^{}}{f_0g\mathrm{sin}\alpha }},`$ (C.11)
$`g_4`$ $`=`$ $`{\displaystyle \frac{\mathrm{cos}\beta \beta ^{}(h^21)}{f_0g(h^2+1)}},`$ (C.12)
$`g_5`$ $`=`$ $`{\displaystyle \frac{\beta ^{}\mathrm{cot}\alpha (h^21)}{f_0g(h^2+1)}},`$ (C.13)
$`g_6`$ $`=`$ $`{\displaystyle \frac{i\beta ^{}(h^21)}{f_0g\mathrm{sin}\alpha (h^2+1)}}.`$ (C.14)
## Appendix D The resolution of an $`A_1`$ singularity
The Eguchi-Hansen metric can be written as
$$ds^2=4r^6(a+r^4)^{\frac{5}{2}}(2a+r^4)dr^2+(a+r^4)^{\frac{1}{2}}\left(\frac{r^8}{2a+r^4}(\sigma ^1)^2+(2a+r^4)((\sigma ^2)^2+(\sigma ^3)^2)\right),$$
(D.1)
with $`r0`$. This is a global coordinate system which allows a smooth $`a0`$ limit. A corresponding vielbein is
$`e^1`$ $`=`$ $`A_1\sigma ^1=r^4(a+r^4)^{\frac{1}{4}}(2a+r^4)^{\frac{1}{2}}\sigma ^1,`$ (D.2)
$`e^2`$ $`=`$ $`A_2\sigma ^2=(a+r^4)^{\frac{1}{4}}(2a+r^4)^{\frac{1}{2}}\sigma ^2,`$ (D.3)
$`e^3`$ $`=`$ $`A_3\sigma ^3=(a+r^4)^{\frac{1}{4}}(2a+r^4)^{\frac{1}{2}}\sigma ^3,`$ (D.4)
$`e^4`$ $`=`$ $`A_4\sigma ^4=2r^3(a+r^4)^{\frac{5}{4}}(2a+r^4)^{\frac{1}{2}}dr.`$ (D.5)
The linearized perturbation around $`a=0`$ is given by
$`\delta A_1`$ $`=`$ $`{\displaystyle \frac{5}{4r^3}}\delta a,`$ (D.6)
$`\delta A_2=\delta A_3`$ $`=`$ $`{\displaystyle \frac{3}{4r^3}}\delta a,`$ (D.7)
$`\delta A_4`$ $`=`$ $`{\displaystyle \frac{1}{2r^4}}\delta a`$ (D.8)
The size of the deformation can be determined using the natural metric
$$\delta g_{ij}\delta g_{kl}g^{ik}g^{jl}=\sqrt{A_1A_2A_3A_4}\underset{i}{}\left(\frac{2\delta A_i}{A_i}\right)^2=\sqrt{2}r^{\frac{13}{2}}\delta a^2.$$
(D.9)
This diverges as $`r0`$, which indicates that the range of validity of the linearized approximation is smaller for small $`r`$.
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# Three Charge Supertubes in Type IIB Plane Wave Backgrounds
## 1 Introduction
Following the discovery of the vacuum rotating black ring<sup>1</sup><sup>1</sup>1Although another rotating black ring has recently been found , due to its conical singularity it is strictly not a vacuum solution. a number of generalisations have been constructed. These could be organised broadly into two groups. The first encompasses asymptotically flat charged rotating rings, which include both the supersymmetric and non-extremal black rings of minimal supergravity as well as the dipole rings found in . The second set consists of non-asymptotically flat solutions. It is relatively straightforward to use solution-generating techniques to construct black rings in fluxbranes . More interesting has been the recent construction of supersymmetric black rings in Taub-NUT backgrounds . Furthermore, Ortin has made use of the fact one can deform supersymmetric solutions of minimal supergravity in order to derive a black ring that asymptotes to the maximally supersymmetric Gödel spacetime . In this note we show how one can trivially generalise this procedure to the more general case of minimal supergravity coupled to $`N1`$ abelian vector multiplets. While these solutions are interesting in their own right, we will be interested in using them as a means to construct supertube configurations in ten and eleven dimensions.
One can lift solutions of the $`U(1)^3`$ five dimensional supergravity to eleven dimensions. When one does this with the supersymmetric black ring solution, one obtains a black supertube. Via dimensional reduction and a series of certain T-dualities, they represent D1-D5-P supertubes in Type IIB supergravity . On the other hand, performing this sequence of dualities on the Gödel solution, one obtains a supersymmetric plane wave . Thus one would expect that lifting Gödel black rings would lead to supertubes embedded in a plane wave. Indeed a three charge Gödel BMPV black hole was constructed in using the reverse procedure. In that work, an asymptotically plane wave D1-D5-P configuration was constructed, which upon T-duality and dimensional reduction led to a three-charge Gödel BMPV black hole. In this note we construct three charge black supertubes that are asymptotically plane wave. We find that the near horizon geometry is unchanged by the deformation. These generalise the supertubes of . We find the presence of the Kaluza-Klein monopole tube somewhat obstructs the inclusion of the plane wave. More precisely, if the Kaluza Klein dipole charge $`q^3`$ is non zero, our solutions necessarily contain closed timelike curves. Further, we can only construct a three charge supertube in the maximally supsersymmetric plane wave in the case where $`q^3=0`$.
This paper is organised as follows. First, we review the construction and supersymmetric solutions of the general five dimensional minimal supergravity coupled to $`N1`$ abelian vector multiplets. We show how these can be deformed simply. Next, we uplift the solution and present the resulting Type IIB configuration. We show explicitly that far from the branes it asymptotes to a supersymmetric plane wave, supported by the appropriate Ramond-Ramond three form flux. In the following section we consider embedding the three charge supertubes in more general supersymmetric plane wave backgrounds. In particular we derive a solution with three charges and two dipole charges that asymptotes to the maximally supersymmetric plane wave. Finally, we conclude with a discussion and comment on extensions of the work.
## 2 Asymptotically Gödel supersymmetric black rings
To begin we will concern ourselves with $`D=5`$ minimal supergravity coupled to $`N1`$ abelian vector multiplets with scalars valued in a symmetric space. We follow the notation of . The action of this theory is:
$`S={\displaystyle \frac{1}{16\pi }}{\displaystyle }(R1G_{IJ}dX^IdX^JG_{IJ}F^IF^J{\displaystyle \frac{1}{6}}C_{IJK}F^IF^JA^K)`$ (1)
where $`I,J,K=1,\mathrm{},N`$ and the scalars $`X^I`$ are constrained by
$$\frac{1}{6}C_{IJK}X^IX^JX^K=1$$
(2)
where $`C_{IJK}=C_{(IJK)}`$ and the following condition is obeyed:
$$C_{IJK}C_{J^{}(LM}C_{PQ)K^{}}\delta ^{JJ^{}}\delta ^{KK^{}}=\frac{4}{3}\delta _{I(L}C_{MPQ)}.$$
(3)
The matrix $`G_{IJ}`$ is defined as
$$G_{IJ}=\frac{9}{2}X_IX_J\frac{1}{2}C_{IJK}X^K$$
(4)
where one lowers indices on the scalars as follows:
$$X_I=\frac{1}{6}C_{IJK}X^JX^K.$$
(5)
The classification of supersymmetric solutions of this theory can be deduced from that of the gauged theory initiated in and subsequently generalized in . Such solutions admit a globally defined non-spacelike Killing vector field $`V`$. If there exists a neighbourhood where $`V`$ is timelike one can choose coordinates $`(t,x^m)`$ such that $`V=/t`$ and
$$ds^2=f^2(dt+\omega )^2+f^1h$$
(6)
where $`h`$ is a Riemmannian metric on the base-space $``$, $`f`$ is a function and $`\omega `$ a 1-form both living on $``$. Let $`e^0=f(dt+\omega )`$ and we will choose an orientation of $``$ such that $`e^0\text{vol}(h)`$ is positively oriented. We decompose $`d\omega `$ into self-dual and anti self-dual parts on the base as
$$fd\omega =G^++G^{}.$$
(7)
Supersymmetry then implies that $`h`$ is a hyper-Kähler metric on $``$ and that
$$F^I=d(X^Ie^0)+\mathrm{\Theta }^I$$
(8)
where $`\mathrm{\Theta }^I`$ are self-dual two forms on $``$ such that
$$X_I\mathrm{\Theta }^I=\frac{2}{3}G^+.$$
(9)
The Bianchi identities for $`F^I`$ then give $`d\mathrm{\Theta }^I=0`$. The Maxwell equations imply
$$^2(f^1X_I)=\frac{1}{6}C_{IJK}\mathrm{\Theta }^J\mathrm{\Theta }^K,$$
(10)
where $`\alpha \beta \frac{1}{p!}\alpha ^{m_1m_2\mathrm{}m_p}\beta _{m_1m_2\mathrm{}m_p}`$ for $`p`$-forms $`\alpha `$ and $`\beta `$ on $``$.
Ortin has made the interesting observation that a solution of the minimal theory can be deformed by adding a piece to $`G^{}`$ while still leaving it a solution. This obviously generalises to the case considered here where we have a minimal theory coupled to $`N1`$ vector multiplets. He showed that by a judicious choice of $`G^{}`$ one can make the supersymmetric ring asymptotically Gödel. We show here that the same construction works for the supersymmetric ring of the $`U(1)^N`$ theory.
Firstly, we write down the $`U(1)^N`$ supersymmetric black ring. The base $`(,h)`$ is $`𝔼^4`$. We write the metric in toroidal coordinates as
$$h=\underset{i=1}{\overset{4}{}}dx_i^2=\frac{R^2}{(xy)^2}\left((y^21)d\psi ^2+\frac{dy^2}{y^21}+(1x^2)d\varphi ^2+\frac{dx^2}{1x^2}\right).$$
(11)
The various quantities are given by
$`\mathrm{\Theta }^I`$ $`=`$ $`{\displaystyle \frac{1}{2}}q^I(dyd\psi +dxd\varphi ),`$ (12)
$`f^3`$ $`=`$ $`{\displaystyle \frac{1}{6}}C^{IJK}H_IH_JH_K,`$ (13)
$`{\displaystyle \frac{1}{3}}H_I`$ $``$ $`f^1X_I=\overline{X}_I+{\displaystyle \frac{1}{6R^2}}\left(Q_I{\displaystyle \frac{1}{2}}C_{IJK}q^Jq^K\right)(xy)`$ (14)
$`{\displaystyle \frac{1}{24R^2}}C_{IJK}q^Jq^K(x^2y^2),`$
$`\omega `$ $`=`$ $`\omega _\varphi d\varphi +\omega _\psi d\psi ,`$ (15)
$`\omega _\varphi `$ $`=`$ $`{\displaystyle \frac{1}{8R^2}}(1x^2)[q^IQ_I\varrho (3+x+y)],`$
$`\omega _\psi `$ $`=`$ $`{\displaystyle \frac{3}{2}}(1+y)q^I\overline{X}_I{\displaystyle \frac{1}{8R^2}}(y^21)[q^IQ_I\varrho (3+x+y)],`$
where $`q^I`$, $`Q_I`$, $`\overline{X}_I`$ are constants, $`C^{IJK}=C_{IJK}`$ and $`\varrho =\frac{1}{6}C_{IJK}q^Iq^Jq^K`$. The constants $`\overline{X}_I`$ obey the same constraint as $`X_I`$ do. The coordinate ranges are, as usual for black rings, $`1x1`$ and $`\mathrm{}<y1`$ and both angles have period $`2\pi `$. There is an event horizon at $`y=\mathrm{}`$. Using the method of Ortin, we deform the solution as follows. Simply replace $`\omega `$ by $`\omega ^{}=\omega +\omega _G`$ where $`\omega _G`$ is given by
$$\omega _G=\frac{\mu R^2}{(xy)^2}[(1x^2)d\varphi (y^21)d\psi ].$$
(16)
It is easy to check that $`(d\omega _G)^+=0`$<sup>2</sup><sup>2</sup>2The orientation is defined by $`ϵ_{y\psi x\varphi }=+1`$ as in and corresponds to $`ϵ_{x_1x_2x_3x_4}=+1`$. and thus $`G^+`$ for this deformed solution is the same as in the undeformed case. This is what implies it is still a solution. The remarkable fact is that this deformation leaves the horizon intact and thus the solution still represents a black ring. This is easy to see since as $`y\mathrm{}`$ the extra terms arising in the metric from $`\omega _G`$ vanish. Also, as promised, the solution asymptotes to the maximally supersymmetric Gödel spacetime. To see this one needs to introduce the polar coordinates
$$\rho \mathrm{sin}\theta =\frac{R\sqrt{y^21}}{xy},\rho \mathrm{cos}\theta =\frac{R\sqrt{1x^2}}{xy}$$
(17)
where $`0\theta \pi /2`$ and $`0\rho <\mathrm{}`$. Then using the fact that the undeformed solution is asymptotically flat it is easy to deduce that as $`\rho \mathrm{}`$
$$ds^2(dt+\omega _G)^2+d\rho ^2+\rho ^2(d\theta ^2+\mathrm{sin}^2\theta d\psi ^2+\mathrm{cos}^2\theta d\varphi ^2).$$
(18)
In these coordinates $`\omega _G=\mu \rho ^2\sigma _R^3/2`$ and $`\sigma _R^3=d\varphi ^{}+\mathrm{cos}\theta ^{}d\psi ^{}`$ is a right invariant form<sup>3</sup><sup>3</sup>3The Euler angles $`(\theta ^{},\varphi ^{},\psi ^{})`$ are given by $`\theta ^{}=2\theta `$, $`\psi ^{}=\varphi +\psi `$ and $`\varphi ^{}=\varphi \psi `$. on the $`S^3`$. This corresponds to the maximally supersymmetric Gödel solution. Note that this particular deformation does not introduce any Dirac-Misner strings ($`\omega _\varphi ^{}(x=\pm 1)=\omega _\psi ^{}(y=1)=0`$) but closed timelike curves (CTC) will of course occur. We should note that the near-horizon limit of this deformed black ring is still locally $`AdS_3\times S^2`$. This is easy to see using the technique of . The extra terms present in the metric due to the deformation are
$$f^2\omega _G^22f^2(dt+\omega )\omega _G$$
(19)
and as $`y\mathrm{}`$, $`f^2\text{const}\times y^4`$, $`\omega \text{const}\times y^3d\psi `$ and $`\omega _G\mu R^2d\psi `$. Then letting<sup>4</sup><sup>4</sup>4The radius of the $`S^1`$ of the ring at the horizon is denoted by $`L`$ and is a function of the charges and $`R`$ as given in . $`\stackrel{~}{r}=R^2/(ϵLy)`$ and $`\stackrel{~}{t}=t/ϵ`$ it is clear that as $`ϵ0`$ both of the above terms vanish. Note that the extra term present in each of the gauge fields also vanishes in this limit.
The solution we have constructed can therefore be seen to describe asymptotically Gödel supersymmetric black rings coupled to $`N1`$ vector multiplets. As in the undeformed case, they are $`\frac{1}{2}`$ BPS. We note in passing that in the $`R=0`$ limit one recovers the analogous supersymmetric Gödel BMPV black hole. This is made manifest by passing into the $`(\rho ,\theta )`$ coordinate system.
## 3 Asymptotically plane wave D1-D5-P supertube
It is well known that $`D=11`$ supergravity, reduced on a $`T^6`$ with coordinates $`z_a`$, $`a=\{1,\mathrm{},6\}`$, using the Ansatz
$`ds_{11}^2`$ $`=`$ $`ds_5^2+X^1(dz_1^2+dz_2^2)+X^2(dz_3^2+dz_4^2)+X^3(dz_5^2+dz_6^2)`$
$`C_3`$ $`=`$ $`A^1dz_1dz_2+A^2dz_3dz_4+A^3dz_5dz_6`$ (20)
yields the STU model. This has the action (1) with $`N=3`$, the only non vanishing component of $`C_{IJK}`$ is $`C_{123}=1`$ and permutations, and the matrix $`G_{IJ}=\frac{1}{2}(X^I)^2\delta _{IJ}`$. In this particular case, the solution $`ds_5^2`$ is given by (6) along with (11)-(15) and the replacement $`\omega \omega ^{}`$, describes a $`\frac{1}{2}`$-BPS three-charge Gödel black ring. Note that the field strengths (8) are also deformed. Lifting these black rings to eleven dimensional supergravity then yields a straightforward deformation of the three-charge M-theory supertubes given in . Explicitly, in terms of the $`H_I`$, we rewrite (6) as
$$ds_5^2=(H_1H_2H_3)^{\frac{2}{3}}(dt+\omega ^{})^2+(H_1H_2H_3)^{\frac{1}{3}}\underset{i=1}{\overset{4}{}}dx_i^2$$
(21)
and the 1-form potentials are
$$A^I=H_I^1(dt+\omega ^{})\frac{q^I}{2}((1+y)d\psi +(1+x)d\varphi ).$$
(22)
Here, as in we have set $`\overline{X}_I=\frac{1}{3}`$. The four supercharges of the five dimensional solution are inherited to yield $`\frac{1}{8}`$ BPS configuration. The undeformed system presented in consists of three M2 branes carrying conserved charges proportional to the $`Q^I`$, and three M5 branes, each of which wrap the ‘ring’ $`\psi `$ coordinate which is transverse to the membranes. As explained clearly in , these M5 branes do not carry conserved charges but instead possess ‘dipole’ charges parameterised by the $`q^I`$. It should be noted that the notion of mass for these objects is defined relative to an asymptotic Minkowskian region transverse to the M2 branes. However, rather than being asymptotically flat, the solution presented here obviously asymptotes to the $`\frac{5}{8}`$ BPS supersymmetric Gödel universe that arises upon lifting (18).
To construct the Type IIB solution, Kaluza-Klein reduce on the compact direction $`z_6`$ and T-dualize along $`z_5,z_4,z_3`$. We could of course choose any of the $`z`$ as the initial $`S^1`$. This choice corresponds to taking $`z_5`$ to point along the axis of the supertube. The resulting string frame metric is
$$ds_{\text{IIB}}^2=\frac{(dt+\omega ^{})^2}{H_3\sqrt{H_1H_2}}+\frac{H_3}{\sqrt{H_1H_2}}(dz_5+A^3)^2+\sqrt{H_1H_2}\underset{i=1}{\overset{4}{}}dx_i^2+\sqrt{\frac{H_2}{H_1}}\underset{i=1}{\overset{4}{}}(dz_i)^2$$
(23)
with dilaton
$$e^{2\mathrm{\Phi }}=\frac{H_2}{H_1}$$
(24)
and three-form RR field strength
$$F_3=(X^1)^2_5F^1+F^2(dz_5+A^3).$$
(25)
The solution above is similar to the D1-D5-P ‘double-helix’ supertube, except now it is not asymptotically flat. Given the fact that Gödel spacetimes are T-dual to plane waves, we expect something similar to manifest itself in this solution. To see this explicitly let us study the asymptotic form of this IIB solution. As $`\rho \mathrm{}`$
$$ds_{\text{IIB}}^2dz_5^2+2dz_5(dt+\omega _G)+\underset{i=1}{\overset{4}{}}dx_i^2+\underset{i=1}{\overset{4}{}}dz_i^2$$
(26)
and if we let $`Z=t+z_5`$ we get
$$ds_{\text{IIB}}^2dt^2+dZ^2+2(dZdt)\omega _G+\underset{i=1}{\overset{4}{}}dx_i^2+\underset{i=1}{\overset{4}{}}dz_i^2.$$
(27)
Note that this form of the metric is just as that found in after T-dualising the Gödel IIA solution. Thus if one makes the coordinate transformation
$`u=tZ,v=t+Z`$ (28)
$`\stackrel{~}{\varphi }=\varphi \mu u,\stackrel{~}{\psi }=\psi +\mu u,`$ (29)
$`\stackrel{~}{x}_1+i\stackrel{~}{x}_2=r_1e^{i\stackrel{~}{\varphi }},\stackrel{~}{x}_3+i\stackrel{~}{x}_4=r_2e^{i\stackrel{~}{\psi }}`$ (30)
and noting $`\rho ^2=r_1^2+r_2^2`$, the asymptotic form of the metric becomes
$$ds_{\text{IIB}}^2dudv\mu ^2\left(\underset{i=1}{\overset{4}{}}\stackrel{~}{x}_i^2\right)du^2+\underset{i=1}{\overset{4}{}}d\stackrel{~}{x}_i^2+\underset{i=1}{\overset{4}{}}dz_i^2$$
(31)
which is a 1/2 supersymmetric plane wave solution to type IIB. The flux becomes
$$F_3\frac{\mu }{2}dud(\rho ^2\stackrel{~}{\sigma _R^3}).$$
(32)
In fact it is the Penrose limit of $`AdS_3\times S^3\times T^4`$ supported by a three form RR-flux, which can be derived as the S-dual of $`AdS_3\times S^3\times T^4`$ supported by an NS-NS three form . Hence the solution (23) seems to represent a D1-D5-P supertube in this plane wave background. One needs to check that this IIB solution has a regular horizon. Since we have noted that $`ds_5^2`$ is regular as $`y\mathrm{}`$ we simply need to show that $`dz_5+A^3`$ is also regular. One can do this in exactly the same manner as was done in by performing the same shift in $`z_5`$ as they did. This is because the extra term we have in $`A^3`$ is $`H_3^1\omega _G`$ which vanishes as $`y\mathrm{}`$. In fact one can go further and show that the near-horizon limit of this IIB solution is unchanged by the deformation and thus is locally $`AdS_3\times S^3\times T^4`$. This can be deduced from the fact that $`ds_5^2`$ has the same near-horizon limit as the undeformed case together with the fact that the extra term in $`A^3`$ is $`O(\overline{r}^2)`$, where $`\overline{r}=R/y`$. Thus the IIB solution we have constructed interpolates between $`AdS_3\times S^3\times T^4`$ and its Penrose limit.
A subtlety concerning the global structure of the spacetime derived here should be noted. Since the Kaluza-Klein direction $`z_5`$ is compact with period $`2\pi R_z`$ as in the undeformed solution, the lightcones coordinates $`u=z_5`$ and $`v=2t+z_5`$ inherit this periodicity. Indeed, enforcing the coordinate transformation (29) to be valid globally one deduces that $`2\pi R_z\mu =2\pi N`$ for some integer $`N`$. Thus we see that the strength of the flux of the plane wave $`\mu =\frac{N}{R_z}`$ is quantized in units of inverse radius of the compact direction $`z_5`$. Asymptotically the deformed solution is actually a discrete quotient of a plane wave. Hence we will have CTC since $`\frac{}{u}`$ is timelike in this plane wave. It is not surprising that we have this situation, given that before the T-dualities we were dealing with an asymptotically Gödel spacetime which contains CTC at every point for sufficiently large radius. However in the IIB solution the CTC are milder in the sense that they are global, just like the ones encountered in AdS spacetimes. Moreover, these features are not specific to our solutions, but also occur for the D1-D5-P systems of . Curiously, one cannot pass to the covering space where $`u`$ is a non-compact coordinate as $`z_5`$ has to be periodic when $`q^3>0`$ in order to avoid a Dirac-Misner string singularity at $`x=1`$ . Thus unlike the BMPV case we cannot remove these CTC, unless $`q^3=0`$ in which case the horizon is singular.
The microscopic derivation of the entropy of the asymptotically flat supertubes has been examined in , though no complete derivation has been given in terms of the D1-D5-P CFT. Moreover there has been success from the eleven dimensional standpoint by considering the CFT of the M5 brane intersection . It is interesting to note that since the geometry of the horizon is unaffected by the deformation we introduced, the entropy of this configuration is the same as that of the asymptotically flat three charge supertubes found in . Now, strings in the Penrose limit of $`AdS_3\times S^3\times T^4`$ can be easily quantized . It turns out that the asymptotic density of states for this model is the same as in flat space (as $`\alpha ^{}0`$). This is nice as it suggests that given the microscopic derivation of the entropy of the three charge supertube in flat space one should find the same answer as for the three charge supertube in the plane wave solution constructed in this paper. This is consistent with the fact that the Bekenstein-Hawking entropy is unaffected by the plane wave as remarked above.
## 4 Supertubes in the maximally supersymmetric plane wave background
The solution we have constructed may be reduced along $`S^1\times T^4`$ to yield the seed Gödel solution we started with. This of course requires the metric to be independent of the compactification coordinates. On the other hand, one may wish to consider supertube configurations embedded in more general plane wave backgrounds, such as the maximally supersymmetric plane wave background.
Consider the undeformed solution. The Einstein frame metric is given in terms of the string frame metric as $`(g_E)_{\mu \nu }=e^{\varphi /2}(g_S)_{\mu \nu }`$. For the case at hand
$`ds_E^2`$ $`=`$ $`H_1^{1/4}H_2^{3/4}\left[2(dt+\omega )(dz_5+\mathrm{\Omega })+H_3(dz_5+\mathrm{\Omega })^2\right]`$ (33)
$`+`$ $`H_1^{3/4}H_2^{1/4}{\displaystyle \underset{i=1}{\overset{4}{}}}dx_i^2+\left({\displaystyle \frac{H_2}{H_1}}\right)^{1/4}{\displaystyle \underset{i=1}{\overset{4}{}}}dz_i^2,`$
$`\mathrm{\Omega }`$ $``$ $`{\displaystyle \frac{q^3}{2}}((1+x)d\varphi +(1+y)d\psi ).`$ (34)
Changing variables<sup>5</sup><sup>5</sup>5Note that in this section our $`u,v`$ coordinates are defined differently to the previous section, in order to match with . to $`2u=z_5`$ and $`2v=2t+z_5`$ transforms the terms in the square bracket to
$$4dudv2du(2\omega +\mathrm{\Omega })+(2dv+2\omega )\mathrm{\Omega }+\mathrm{\Omega }^2+(H_31)(2du\mathrm{\Omega })^2.$$
(35)
This can be cast in a nicer form by introducing a new $`u`$ coordinate by $`dudu+\mathrm{\Omega }/2`$ which leaves the full metric as
$`ds_E^2`$ $`=`$ $`H_1^{1/4}H_2^{3/4}\left[4dudv2du(2\omega +\mathrm{\Omega })+4(H_31)du^2\right]`$ (36)
$`+`$ $`H_1^{3/4}H_2^{1/4}{\displaystyle \underset{i=1}{\overset{4}{}}}dx_i^2+\left({\displaystyle \frac{H_2}{H_1}}\right)^{1/4}{\displaystyle \underset{i=1}{\overset{4}{}}}dz_i^2.`$
This is rather nice as we have cast the D1-D5-P supertube in exactly the same form as the D1-D5-P system which reduces to the BMPV black hole. This form of the metric allows one to add in extra pieces to the $`g_{uu}`$ component very easily as was done for the BMPV system in . This relies on the following observation:
$`ds^2`$ $`=`$ $`e^{2A}\left(4dudv+du^2+du{\displaystyle \underset{i=1}{\overset{n}{}}}C_idx_i\right)+{\displaystyle \underset{i=1}{\overset{n}{}}}e^{2B_i}dx_i^2`$ (37)
$`R_{\mu \nu }`$ $`=`$ $`\overline{R}_{\mu \nu }{\displaystyle \frac{1}{2}}\delta _\mu ^u\delta _\nu ^ue^{2A}{\displaystyle \underset{i=1}{\overset{n}{}}}e^{2B_i}[_i^2+_i_iG_i+2(_i^2A+_iA_iG_i)]`$ (38)
$`G_i`$ $`=`$ $`2A2B_i+{\displaystyle \underset{j=1}{\overset{n}{}}}B_j`$ (39)
where $`A,B_i,C_i`$ are all functions of the transverse coordinates $`x_i`$ and $`\overline{R}_{\mu \nu }`$ is the Ricci tensor with $`=0`$<sup>6</sup><sup>6</sup>6In it was stated that the inclusion of angular momentum does not affect the result quoted in which had $`C_i=0`$..
Armed with these results we now deform the D1-D5-P supertube as follows:
$`ds_E^2`$ $`=`$ $`H_1^{1/4}H_2^{3/4}\left[4dudv2du(2\omega +\mathrm{\Omega })+(+4(H_31))du^2\right]`$ (40)
$`+`$ $`H_1^{3/4}H_2^{1/4}{\displaystyle \underset{i=1}{\overset{4}{}}}dx_i^2+\left({\displaystyle \frac{H_2}{H_1}}\right)^{1/4}{\displaystyle \underset{i=1}{\overset{4}{}}}dz_i^2.`$
In view of the general result quoted above it is immediate that the the Ricci tensor of this deformed configuration is:
$`R_{\mu \nu }=\overline{R}_{\mu \nu }{\displaystyle \frac{1}{2H_1H_2}}\delta _\mu ^u\delta _\nu ^u\left(_{x_i}^2+H_1_{z_i}^2{\displaystyle \frac{1}{4}}\left(_{x_i}^2\mathrm{log}H_1+3_{x_i}^2\mathrm{log}H_2\right)\right).`$ (41)
It should be noted that in this case the functions $`G_i=0`$ and thus there are no terms with first derivatives of $``$. We will choose to support this deformation by a five form flux as in . Note that this is in contrast to the previous section where the deformation was supported by a three form flux. The five form flux must be self dual and $`F_5F_3=0`$. It is natural to try and use the same expression as in . Namely
$$F_5=\mu du(dx_1dx_2dz_1dz_2+dx_3dx_4dz_3dz_4).$$
(42)
It is in fact straightforward to check that this five form is self dual also in this case. However in general $`F_3F_50`$, see Appendix. In the special case $`q^3=0`$ though $`F_3F_5=0`$ and thus we may straightforwardly deform the solution in this situation. Note that the supertube is nakedly singular in this limit though. Nevertheless as noted in the world volume theory of this configuration has a sensible interpretation in Type IIA in terms of D6 branes.
The analysis follows similarly to . Explicitly, the Type IIB Einstein equations are
$$R_{\mu \nu }=\frac{1}{2}_\mu \varphi _\nu \varphi +\frac{1}{96}F_{\mu \alpha \beta \gamma \delta }F_\nu ^{\alpha \beta \gamma \delta }+\frac{e^\varphi }{4}\left(F_{\mu \alpha \beta }F_\nu ^{\alpha \beta }\frac{1}{12}g_{\mu \nu }F_3^2\right)$$
(43)
Clearly only the $`uu`$ component of the stress energy tensor is altered by the presence of $``$. It is not difficult to check that these extra terms are
$`\mathrm{\Delta }(F_{\mu \alpha \beta }F_\nu ^{\alpha \beta })`$ $`=`$ $`2(_i\mathrm{log}H_2)^2H_1^{\frac{1}{2}}H_2^{\frac{3}{2}}\delta _\mu ^u\delta _\nu ^u`$ (44)
$`\mathrm{\Delta }(g_{\mu \nu }F_3^2)`$ $`=`$ $`6((_i\mathrm{log}H_1)^2(_i\mathrm{log}H_2)^2)H_1^{\frac{1}{2}}H_2^{\frac{3}{2}}\delta _\mu ^u\delta _\nu ^u`$ (45)
$`F_{\mu \alpha \beta \gamma \delta }F_\nu ^{\alpha \beta \gamma \delta }`$ $`=`$ $`{\displaystyle \frac{48\mu ^2}{H_1H_2}}\delta _\mu ^u\delta _\nu ^u`$ (46)
where $`\mathrm{\Delta }`$ represents the change in the quantity in the brackets due to the deformation $``$. The equations of motion for the dilaton and $`F_3`$ are unchanged. Enforcing the $`uu`$ component of the Einstein equations then leads to a simple linear equation for the deformation $``$:
$$(_{x_i}^2+H_1_{z_i}^2)=\mu ^2.$$
(47)
Note that to do this one needs to use the fact that $`H_1,H_2`$ are harmonic when $`q^3=0`$. Thus the equation for the deformation is identical to that found in for the BMPV system, except now $`H_1`$ is a harmonic function with delta function sources on a ring $`\rho =R`$ and $`\theta =\pi /2`$ as opposed to at the origin. This makes the PDE (47) more complicated and in fact if one writes it in $`(\rho ,\theta )`$ or $`(x,y)`$ coordinates it is not separable. Remarkably, if one uses yet a different coordinate system the equation can be made separable. The coordinates in question are:
$$r^2=R^2\frac{1x}{xy},\mathrm{cos}^2\mathrm{\Theta }=\frac{1+x}{xy}$$
(48)
and were introduced in . The flat space metric then looks like
$$\underset{i=1}{\overset{4}{}}dx_i^2=\mathrm{\Sigma }\left(\frac{dr^2}{r^2+R^2}+d\mathrm{\Theta }^2\right)+(r^2+R^2)\mathrm{sin}^2\mathrm{\Theta }d\psi ^2+r^2\mathrm{cos}^2\mathrm{\Theta }d\varphi ^2$$
(49)
where $`\mathrm{\Sigma }=r^2+R^2\mathrm{cos}^2\mathrm{\Theta }`$. It is immediately apparent that there are solutions of the form $`=X(x_i)+Z(z_i)`$ where $`_{z_i}^2Z=\alpha ^2`$ where $`\alpha ^2`$ is a separation constant. The resulting equation for $`X`$ is then $`_{x_i}^2X+\alpha ^2H_1=\mu ^2`$, which upon multiplication by $`\mathrm{\Sigma }`$ is additively separable in the $`r,\mathrm{\Theta }`$ coordinates system since $`H_1=1+Q_1/\mathrm{\Sigma }`$. This means that $`X=F(r)+G(\mathrm{\Theta })`$. The function $`F(r)`$ satisfies
$`{\displaystyle \frac{1}{r}}{\displaystyle \frac{d}{dr}}\left(r(r^2+R^2){\displaystyle \frac{dF}{dr}}\right)+r^2(\alpha ^2+\mu ^2)`$ $`=`$ $`\beta ^2`$ (50)
where $`\beta ^2`$ is another separation constant. This may be be integrated to give
$`F(r)`$ $`=`$ $`{\displaystyle \frac{r^2}{8}}(\alpha ^2+\mu ^2)+{\displaystyle \frac{R^2}{8}}(\alpha ^2+\mu ^2)\mathrm{log}(r^2+R^2){\displaystyle \frac{\beta ^2}{4}}\mathrm{log}(r^2+R^2)`$ (51)
$`+`$ $`c_1\mathrm{log}(r/\sqrt{r^2+R^2})+c_2`$
where $`c_1,c_2`$ are integration constants. The equation for $`G`$ may also be integrated, but it is convenient to change variables to $`z=\mathrm{sin}^2\mathrm{\Theta }`$ first. In terms of $`z`$ we have
$$4\frac{d}{dz}\left(z(1z)\frac{dG}{dz}\right)+R^2(\alpha ^2+\mu ^2)(1z)+Q_1\alpha ^2=\beta ^2$$
(52)
which leads to
$`G(z)`$ $`=`$ $`{\displaystyle \frac{(Q_1\alpha ^2\beta ^2)}{4}}\mathrm{log}(1z)+{\displaystyle \frac{R^2(\alpha ^2+\mu ^2)}{8}}(\mathrm{log}(1z)z)`$ (53)
$`+`$ $`c_3(\mathrm{log}z\mathrm{log}(1z))+c_4.`$ (54)
We have generated quite a few constants upon integrating (47), however they may all be fixed as follows. The constants $`c_2,c_4`$ can be absorbed into shifts of $`v`$. Demanding regularity<sup>7</sup><sup>7</sup>7The horizon in these coordinates is located at $`r=0`$ and $`\mathrm{\Theta }=\pi /2`$, however as already noted it is not regular in the undeformed solution since $`q^3=0`$. at $`\mathrm{\Theta }=0`$ and $`\pi /2`$ forces $`c_3=0`$ and $`\beta ^2=Q_1\alpha ^2+R^2(\alpha ^2+\mu ^2)/2`$ respectively. Demanding regularity at $`r=0`$ forces $`c_1=0`$ and finally requiring that $`\text{const}(x_ix_i+z_iz_i)`$ tells us that $`\alpha ^2=\mu ^2/2`$. Thus we arrive at
$$=\frac{\mu ^2}{16}(r^2+R^2\mathrm{sin}^2\mathrm{\Theta }+z_iz_i)+\frac{1}{8}Q_1\mu ^2\mathrm{log}(r^2+R^2).$$
(55)
This deformation gives a three charge, two dipole supertube which asymptotes to the maximally supersymmetric plane wave solution of type IIB supergravity. The remarks concerning the global causal structure of these deformed supertube configurations given in the previous section also apply here, however in this case we do not get the quantization condition on $`\mu `$. In particular, we can pass to the covering space where $`u`$ is non-compact. One would expect in this case the solution to be devoid of CTC (given the constraint on the charges in ), though we have not performed a full analysis in the intermediate region (in between the near horizon and the asymptotic plane wave).
Finally we note in passing there is a class of solutions (with $`\alpha =0`$) that are independent of the toroidal directions $`z_i`$; in this case, one could easily compactify on $`S^1\times T^4`$ to derive nakedly singular asymptotically Gödel spacetimes in five dimensional minimal supegravity coupled to two vector multiplets.
## 5 Concluding remarks
We have demonstrated how three charge supertubes can be embedded not only in asymptotically flat backgrounds, but also in the next simplest class of solutions, supersymmetric plane wave spacetimes. This was shown explicitly in two cases: firstly, for the Penrose limit of $`AdS_3\times S^3\times T^4`$, and secondly for the maximally supersymmetric BFHP plane wave. The former case was derived by first deforming the supersymmetric black rings in D=5 minimal supergravity coupled to $`N1`$ vector multiplets. In the case $`N=3`$, this solution was uplifted to eleven dimensions to describe three charge supertubes in the supersymmetric Gödel background. Upon Kaluza Klein reduction and T dualisation, a D1-D5-P supertube embedded in $`\frac{1}{2}`$ BPS plane wave was constructed. Furthermore, the inclusion the plane wave term does not affect the properties of the horizon. The removal of Dirac-Misner string singularities, however, leads to CTC, in contrast to the BMPV case.
In the second case, we use a more direct approach to derive a supertube plane wave configuration supported by a self dual five form flux. In order to satisfy the equations of motion, it seems one has to turn off the Kaluza-Klein dipole charge. By noting that the Ricci tensor has a simple decomposition under wave-like deformations of the metric we arrive at a simple linear PDE for the deformation. We show that this is additively separable, in suitable coordinates, and under this condition derive the general solution. A particular solution corresponds to supertubes in the maximally supersymmetric background.
There remain several open problems concerning asymptotically plane wave supertubes. Obviously, one might be interested in their world volume description. This could be particularly relevant for the solutions describing supertubes in the maximally supersymmetric background, as they seem to be nakedly singular from their supergravity description. Furthermore, one could try to generalise the solution presented here in the maximally supersymmetric plane wave background such that one has a non-zero Kaluza Klein dipole charge. Since these three charge, two dipole supertubes have non-extremal counterparts with regular horizons, it might be useful to consider plane wave extensions of these thermally excited supertubes first.
Finally we should emphasise that it is most interesting that one can embed such supertube configurations in non-trivial backgrounds such as plane waves and Gödel spacetimes. These arose from lifting Gödel black rings. There remains, however, the interesting and apparently difficult problem of finding asymptotically AdS black rings.
###### Acknowledgments.
We would like to thank Roberto Emparan for useful comments and reading through a draft of the paper. HKK would like to thank St. John’s College, Cambridge, for financial support.
## Appendix A The three form $`F_3`$
The two form field strengths of the $`U(1)^N`$ five dimensional supergravity are given by:
$$F^I=dA^I=d(X^Ie^0)+\mathrm{\Theta }^I=f^1dH_I^1e^0+H_I^1d\omega +\mathrm{\Theta }^I.$$
(56)
It then follows that the hodge dual is
$$_5F^I=\frac{1}{3!}f^2_iH_I^1ϵ_{ilmn}dx^ldx^mdx^n+H_1^1f^1e^0(G^+G^{})+e^0\mathrm{\Theta }^I$$
(57)
where we have used the fact that $`_4d\omega =f^1(G^+G^{})`$ and $`_4\mathrm{\Theta }^I=\mathrm{\Theta }^I`$. Note that $`ϵ_{0x_1x_2x_3x_4}=1`$ has been used in these Cartesian coordinates. The three form in IIB is given by
$$F_3=(X^1)^2_5F^1+F^2(dz_5+A^3)=(X^1)^2_5F^12F^2du+F^2(A^3\mathrm{\Omega }),$$
(58)
where we have used the new $`u`$ coordinate defined in section 4. Thus we will also need
$`F^2(A^3\mathrm{\Omega })=f^1H_2^1H_3^1d\omega e^0+f^1H_3^1\mathrm{\Theta }^2e^0`$ (59)
$`={\displaystyle \frac{H_1^{1/3}}{H_2^{2/3}H_3^{2/3}}}d\omega e^0+{\displaystyle \frac{H_1^{1/3}H_2^{1/3}}{H_3^{2/3}}}\mathrm{\Theta }^2e^0`$ (60)
and
$`(X^1)^2_5F^1={\displaystyle \frac{1}{3!}}f^2(X^1)^2_iH_1^1ϵ_{lmni}dx^ldx^mdx^n`$
$`+{\displaystyle \frac{H_1^{1/3}}{H_2^{2/3}H_3^{2/3}}}e^0f^1(G^+G^{})+(X^1)^2e^0\mathrm{\Theta }^1.`$ (61)
Upon adding A.5 to A.6 we see that the terms that look like $`e^0G^{}`$ cancel. Interestingly we find further cancellations. Using the identity $`G^+=\frac{1}{2}fH_I\mathrm{\Theta }^I`$ leads to
$`{\displaystyle \frac{H_1^{1/3}}{H_2^{2/3}H_3^{2/3}}}2e^0f^1G^++(X^1)^2e^0\mathrm{\Theta }^1+f^1H_3^1e^0\mathrm{\Theta }^2=f^1H_2^1e^0\mathrm{\Theta }^3.`$ (62)
Putting everything together the final explicit expression for $`F_3`$ is:
$`F_3=2F^2du+{\displaystyle \frac{1}{3!}}f^2(X^1)^2_iH_1^1ϵ_{ilmn}dx^ldx^mdx^nf^1H_2^2e^0\mathrm{\Theta }^3.`$ (63)
Thus as compared to the D1-D5-P BMPV system of we have an extra term $`f^1H_2^2e^0\mathrm{\Theta }^3`$. Note that this term is proportional to $`q^3`$ the KK dipole charge. Thus, shutting off the KK dipole charge of the supertube gets rid of this extra term leaving $`F_3F_5=0`$ as required.
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# Propagation OF Ultra-High Energy Cosmic Rays above 10¹⁹ eV in a structured extragalactic magnetic field and Galactic magnetic field
## 1 INTRODUCTION
The nature of Ultra-High Energy Cosmic Rays (UHECRs), which are particles of energy above $`10^{19}`$ eV, is poorly known. This is one of the most challenging problems of modern astrophysics.
One of problems about UHECRs is what their origin is. The two scenarios of their origin are suggested, which are called bottom-up and top-down ones. On the one hand, bottom-up scenarios assume some astrophysical phenomena as their origin. UHECRs are thought to be of extragalactic origin since the gyroradii of UHECRs above $`10^{19}`$ eV exceed thickness of our galaxy. From this fact and the Hillas plot (Hillas, 1984), probable candidates of UHECR origins are active galactic nuclei (AGNs), gamma-ray bursts (GRBs) and colliding galaxies. Theoretically, this scenario predicts the GZK cutoff of the energy spectrum of UHECRs (Greisen, 1966; Zatsepin & Kuz’min, 1966) since their source candidates are located at far distances. UHE protons with energy above $`4\times 10^{19}`$ eV interact with the cosmic microwave background (CMB) and lose large fraction ($`20\%`$) of their energy per interaction by photopion production (Berezinsky & Grigorieva, 1988; Yoshida & Teshima, 1993). The mean free path of UHE protons through the CMB field is $``$ 10 Mpc at $`10^{20}`$ eV. Thus the energy spectrum at the Earth should have a cutoff around $`E8\times 10^{19}`$ eV. This spectral cutoff is called the GZK cutoff. But there is a observational disagreement of the energy spectra between the Akeno Giant Air Shower Array (AGASA), which does not observe the GZK cutoff (Takeda et al., 1998), and the High Resolution Fly’s Eye (HiRes; Wilkinson et al., 1999), which does it (Abu-Zayyad et al., 2004). This discrepancy between the two experiments remains being one of open questions in astroparticle physics. On the other hand, top-down scenarios assume some processes based on new physics beyond the standard model of the particle physics (see a review Bhattacharjee & Sigl (2000)).
Another problem is arrival distribution of UHECRs. The AGASA reported that there is no statistically significant large scale anisotropy in the observed arrival distribution of UHECRs above $`10^{19}`$ eV (Takeda et al., 1999). This fact points out that sources of UHECRs are distributed isotropically, but isotropic distribution of sources cannot reproduce the small-scale anisotropy reported by the AGASA (Takeda et al., 1999, 2001). A model of UHECR origin are constrained by their ability to reproduce such observed arrival distribution of UHECRs. On the other hand, the HiRes experiment indicates that there is no statistically significant small-scale anisotropy (Abbasi et al., 2004; Farrar et al., 2004). However Yoshiguchi et al. (2004) concluded this discrepancy between the two observations is not statistically significant at present. This problem is left for future investigation by new experiments such as the Pierre Auger Observatory.
To obtain information on UHECR origin, we need to calculate their arrival distribution using some kinds of source models. To do so, we have to simulate propagation of UHECRs in the intergalactic space, where the Extragalactic Magnetic Field (EGMF) plays important roles since we assume that UHECRs are protons in this paper. Yoshiguchi et al. (2003a) calculated propagation of UHE protons in an uniform turbulence of magnetic field with the Kolmogorov spectrum. But such magnetic field is not realistic since the EGMF is expected to reflect the large scale structure of the universe.
In recent years, several groups started to develop physically more realistic models of the EGMF based on numerical simulations of large scale structure formation. Sigl, Miniati, & Ensslin (2003, 2004) used a structured EGMF model which is generated by their large scale structure simulations, taking magnetic fields into account. But their model does not reproduce the local structures actually observed around the Milky Way. This causes the ambiguity in the choice of observer position. In addition to this, the calculated arrival distribution of UHECRs does not correspond to the one expected at the earth.
An important step on modeling the magnetic structure of the local universe is performed by Dolag et al. (2005). They constrain the initial conditions for the dark matter density fluctuations to reproduce the local structures. This allows us to remove the ambiguity in the choice of observer position, and to obtain the simulated skymaps of expected UHE proton deflections in the magnetic large-scale structure around our galaxy. However, they did not calculate the arrival distribution of UHECRs, and thus could not obtain the information on source distribution which reproduces the AGASA observation. And also, the effects of the Galactic Magnetic Field (GMF) are not considered in both Sigl, Miniati, & Ensslin (2003, 2004) and Dolag et al. (2005). But recently it has also been shown that the GMF affects the arrival distribution of UHECRs (Alvarez-Muniz, Engel & Stanev, 2002; Yoshiguchi et al., 2003b). Thus we cannot neglect modifications of arrival directions of UHECRs by the GMF when we simulate their arrival distribution.
In this work, we calculate propagation of UHECRs taking both the EGMF and the GMF into account, and simulate the arrival distribution of UHECRs. We constrain source number density of UHECRs by comparison of the results with the AGASA observation. We generate the magnetic structure of the local universe by our original method (section 3.1) from the IRAS PSCz catalogue of galaxies (Saunders et al., 2000). We also construct our source models of UHECRs from this catalogue. As our GMF model, we adopt a bisymmetric spiral field (BSS) model (Alvarez-Muniz, Engel & Stanev, 2002) just like Yoshiguchi et al. (2003b).
In order to simulate the arrival distribution of UHECRs, we apply a method developed in previous works. (Fl$`\ddot{\mathrm{u}}`$ckiger et al., 1991; Bieber, Evenson, & Lin, 1992; Stanev, 1997; Medina-Tanco, 1999; Yoshiguchi et al., 2003b). We numerically calculate an inverse process of propagation of UHE protons, which reach the earth, and record their trajectories in our Galaxy and the intergalactic space. In other words, we inject UHECRs from the earth isotropically whose charges are taken as -1. We then select some of them according to a given source distribution. (Detailed explanation is given in the section 4.3) The expected arrival distribution can be obtained by mapping the velocity directions of the selected trajectories at the earth. The validity of this method is supported by the Liouville’s theorem. This method enables us to save the CPU time effectively since we calculate only trajectories of UHE protons which reach the earth. A method for this process is explained in section 4.
The outline of this paper is as follows. In section 2, we explain the IRAS PSCz Catalogue and construct our sample of galaxies. In section 3, we introduce our model of the EGMF and the GMF. We explain our numerical methods for calculating arrival distribution of UHECRs and statistical quantities in section 4. Then in section 5, we estimate the most appropriate number density of source of UHECRs and compare the statistical quantities calculated from this source model with the EGMF to those calculated without the EGMF. We also demonstrate skymaps of the arrival distribution of UHECRs. In section 6, we summarize the main results.
## 2 A SAMPLE OF GALAXIES
In order to calculate propagation of UHECRs considering the local structures actually observed around the Milky Way, we use the IRAS PSCz Catalogue (Saunders et al., 2000) for construction of our EGMF model and UHECR source models.
We used the ORS sample of galaxies to construct our UHECR source models in our previous work (Yoshiguchi et al., 2003a). The ORS sample has better completeness on nearby galaxies than the IRAS PSCz catalogue. Thus we used the ORS sample to construct our source model since we were interested in the local sources in order to research nature of UHECRs above $`4\times 10^{19}`$ eV. As a result, we obtained that UHECR source number density is $`10^6\mathrm{Mpc}^3`$. Then we investigated that this source model also explained arrival distribution of UHECRs above $`10^{19}`$ eV though the EGMF was neglected (Yoshiguchi et al., 2003b). In this work, we use galaxy sample to construct not only a source model of UHECRs but also a model of the EGMF, which reflect the large scale structure of the universe. This requests large sky coverage of the galaxy sample. Thus we adopt not the ORS galaxy catalogue but the IRAS PSCz Catalogue.
The IRAS PSCz catalogue consists of 14677 galaxies with redshift and infrared fluxes $`>`$ 0.6 Jy, and covers about 84% of the sky. We assume de Sitter universe with $`\mathrm{\Omega }_m=0.3,\mathrm{\Omega }_\lambda =0.7,H_0=71`$ $`\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ in order to calculate distance of each galaxy. We show distribution of the IRAS PSCz galaxies in Figure 1.
However, there are two problems when we use the IRAS PSCz Catalogue. One is that it is impossible to observe dark galaxies at far-off distance. (the selection effect, see figure 2). The other is that this catalogue has the zone of avoidance (the mask), where the IRAS PSCz Survey did not observe galaxies. Our previous works (Yoshiguchi et al., 2003a, b), which uses the ORS sample (Santiago et al., 1995), also have these problems. Using luminosity function of the IRAS galaxies, we correct these absence of galaxies in the same manner with our previous studies.
We use the luminosity function of the IRAS PSCz galaxies (Takeuchi et al., 2003),
$$\mathrm{\Phi }(L)=\varphi _{}\left(\frac{L}{L_{}}\right)^{1\alpha }\mathrm{exp}\left\{\frac{1}{2\sigma ^2}\left[\mathrm{log}\left(1+\frac{L}{L_{}}\right)\right]^2\right\}.$$
(1)
Here $`L_{}=(4.34\pm 0.86)\times 10^8h^2[L_{}],\alpha =1.23\pm 0.04,\sigma =0.724\pm 0.010,\varphi _{}=(2.34\pm 0.30)\times 10^2h^3\mathrm{Mpc}^3`$. Using this luminosity function, we define the selection function as,
$$\varphi (r)=\frac{_{L_{\mathrm{min}(r)}}^{\mathrm{}}𝑑L\mathrm{\Phi }(L)}{_0^{\mathrm{}}𝑑L\mathrm{\Phi }(L)},$$
(2)
where $`L_{\mathrm{min}(r)}`$ is minimum luminosity of galaxies which are observable at a distance $`r`$. Therefore, $`\varphi (r)`$ represents fraction of all galaxies that are observable at each distance.
First of all, we correct the selection effect. For each of the IRAS galaxies, we add galaxies that are not included in the IRAS sample. The number of added galaxies can be obtained by using the selection function. The positions of added galaxies are determined according to the Gaussian distribution whose mean is the position of the original IRAS galaxy and whose root mean square is $`l(r)`$. Here we define $`l(r)`$ as a mean distance between the original IRAS galaxies at distance $`r`$, thus
$$\frac{4\pi }{3}l(r)^3n(r)=1,$$
(3)
where $`n(r)`$ is number density of the original IRAS galaxies. Then,
$$l(r)=\left(\frac{3}{4\pi }\right)^{1/3}n(r)^{1/3}.$$
(4)
The luminosities of added galaxies are randomly assigned so that their distribution of luminosity is consistent with the luminosity function. This method can complement the IRAS galaxies without spoiling the observed structure of galaxy distribution.
Next, we add galaxies in the mask. We assume that galaxy distribution in this region is homogeneous and number density of galaxies is $`n(r)\varphi (r)^1`$. These luminosities are random but their distribution is consistent with the luminosity function. Our galaxy sample after these corrections is shown in Figure 2.
In this work, we use only the IRAS galaxies within $`100`$ Mpc. We assume that source distribution at $`r>100`$ Mpc is isotropic and uniform, and that their number density is equal to that within $`100`$ Mpc. We neglect cosmological evolution of number density of galaxies. Thus our sample of galaxies reflects the observed local structure within 100 Mpc. We use this sample of galaxies to construct a model of EGMF and source model of UHECRs.
## 3 A MODEL OF MAGNETIC FIELD
### 3.1 Extragalactic Magnetic Field
The EGMF are little known theoretically and observationally. Theoretically, several large scale structure simulations with magnetic field have been performed (Dolag et al., 2005; Sigl, Miniati, & Ensslin, 2004). Roughly speaking, their results are that the strength of magnetic field traces baryon density. A model of Sigl, Miniati, & Ensslin (2004) do not reflect the local structures actually observed around the Milky Way, while one of Dolag et al. (2005) reflects these local structures. To compare model predictions of UHECR arrival distribution with the observed one, it is important to generate the magnetic structure around our galaxy. Therefore, we present a model of the EGMF reflecting the local structures of the universe well.
The EGMF mainly exists in clusters of galaxy or around galaxies. From this standpoint, we present a realistic model of the EGMF. We assume that magnetic field results from the amplification of weak seed fields. In a simulation of evolution of cluster of galaxy (Dolag et al., 2002), the average magnetic field strength in the cluster is amplified as expected from compression alone ($`|B|\rho ^{2/3}`$, where $`\rho `$ is density of matter). We adopt this conclusion
$$|B|\rho ^{2/3}\rho _{L}^{}{}_{}{}^{2/3},$$
(5)
where $`\rho _L`$ is luminosity density explained in the next paragraph. In equation 5, we assume that the luminosity density on each position is proportional to density of the gas. With these assumption, we construct a model of the EGMF as follows.
First, we cover the universe with cubes of side $`l_c`$, which is correlation length of the EGMF. We adopt $`l_c=1\mathrm{M}\mathrm{p}\mathrm{c}`$ from our previous work (Yoshiguchi et al., 2003a). Second, we sum luminosities of galaxies in our sample which exist in each cube. We call this summed value luminosity density. Finally, It is assumed that the magnetic field in each cube is represented as the Gaussian random field with zero mean and has a power-law spectrum
$`B(\stackrel{}{k})B^{}(\stackrel{}{k})`$ $``$ $`k^{n_H}\mathrm{for}2\pi /l_ck2\pi /l_{\mathrm{cut}},`$
$`B(\stackrel{}{k})B^{}(\stackrel{}{k})`$ $`=`$ $`0\mathrm{otherwise},`$ (6)
where $`l_{\mathrm{cut}}`$ is a numerical cutoff scale. Physically, one expects $`l_{\mathrm{cut}}l_\mathrm{c}`$, but we set $`l_{\mathrm{cut}}=1/8\times l_\mathrm{c}`$ in order to save the CPU time. We use $`n_\mathrm{H}=11/3`$, corresponding to the Kolmogorov spectrum since the Faraday rotation map reveals that the clusters’ magnetic fields are turbulent with the Kolmogorov spectrum over at least one order of magnitude of the wavenumber(Vogt & Ensslin, 2004).
Next we consider a normalization of the EGMF. Most observations suggest that clusters of galaxy have magnetic field whose strength is from 0.1 $`\mu `$G to a few $`\mu `$G (see a review Vallee (2004)). On the one hand, the Faraday rotation measurements of polarized radio sources placed within cluster of galaxies provide some evidence for the presence of stronger intracluster magnetic field (ICMF), in the range of a few $`\mu `$G (Taylor et al., 2001; Vogt et al., 2003). On the other hand, observations of hard X-ray emission from cluster of galaxy implies that an average ICMF strength within the emitting volume is 0.2-0.4 $`\mu `$G (Fusco-Femiano et al., 1999; Rephaeli et al., 1999).
In this work, we now set a normalization of its strength $`0.4\mu `$G in a cube where is the center of the Virgo cluster. In order to compare our model with Dolag et al. (2005), we show magnetic field strength within $`100`$ Mpc along three fiducial lines through the Virgo cluster, the Perseus cluster and the Coma cluster in figure 3.1. In figure 3, we also show deflection maps when protons propagate through our EGMF from the distance of 100 Mpc. The deflection angle by uniform turbulent magnetic field is given by
$$\theta 0.3^{}\left(\frac{E}{10^{20}\mathrm{eV}}\right)^1\left(\frac{r}{100\mathrm{M}\mathrm{p}\mathrm{c}}\right)^{1/2}\left(\frac{l_c}{1\mathrm{M}\mathrm{p}\mathrm{c}}\right)^{1/2}\left(\frac{B}{10^4\mu \mathrm{G}}\right).$$
(7)
Here $`r`$ is distance of propagation. In the way similar to Dolag et al. (2005), we assume that proton trajectory makes a random walk through each cell since our EGMF in each cube of side $`l_c`$ is the turbulence. These maps can be compared to figure 13 and 14 in Dolag et al. (2005).
We use our sample of galaxies to construct a model of the EGMF only within $`100`$ Mpc. At a distance above $`100`$ Mpc, we treat the EGMF as an uniform turbulence of magnetic field with the same spectrum and $`|B|`$ = $`1`$ nG since the IRAS PSCz catalogue is poorly covered at $`r>100`$ Mpc.
### 3.2 Galactic Magnetic Field
In this study, we adopt the GMF model used by Alvarez-Muniz, Engel & Stanev (2002), which is composed of the spiral and the dipole field. We briefly explain this GMF model below.
Faraday rotation measurements indicate that the GMF in the disk of the Galaxy has a spiral structure with field reversals at the optical Galactic arms (Beck, 2001). We use a bisymmetric spiral field (BSS) model, which is favored from recent work (Han, Manchester, & Qiao, 1999; Han, 2001). The Solar System is located at a distance $`r_{||}=R_{}=8.5`$ kpc from the center of the Galaxy in the Galactic plane. The local regular magnetic field in the vicinity of the Solar System is assumed to be $`B_{\mathrm{Solar}}1.5\mu \mathrm{G}`$ in the direction $`l=90^\mathrm{o}+p`$ where the pitch angle is $`p=10^\mathrm{o}`$ (Han & Qiao, 1994).
In the polar coordinates $`(r_{||},\varphi )`$, the strength of the spiral field in the Galactic plane is given by
$$B(r_{||},\varphi )=B_0\left(\frac{R_{}}{r_{||}}\right)\mathrm{cos}\left(\varphi \beta \mathrm{ln}\frac{r_{||}}{r_0}\right),$$
(8)
where $`B_0=4.4\mu `$G, $`r_0=10.55`$ kpc and $`\beta =1/\mathrm{tan}p=5.67`$. The field decreases with Galactocentric distance as $`1/r_{||}`$ and it becomes zero for $`r_{||}>20`$ kpc. In the region around the Galactic center ($`r_{||}<4`$ kpc) the field is highly uncertain, and thus assumed to be constant and equal to its value at $`r_{||}=4`$ kpc.
The spiral field strengths above and below the Galactic plane are taken to decrease exponentially with two scale heights (Stanev, 1997),
$$|B(r_{||},\varphi ,z)|=|B(r_{||},\varphi )|\{\begin{array}{ccc}\mathrm{exp}(|z|):\hfill & |z|0.5\mathrm{kpc}& \\ \mathrm{exp}(\frac{3}{8})\mathrm{exp}(\frac{|z|}{4}):\hfill & |z|>0.5\mathrm{kpc}& \end{array}$$
(9)
where the factor $`\mathrm{exp}(3/8)`$ makes the field continuous on $`z`$. The BSS spiral field we use is of even parity, that is, the field direction is preserved at disk crossing.
Observations show that the field in the Galactic halo is much weaker than that in the disk. In this work we assume that the regular field corresponds to a A0 dipole field as suggested in (Han, 2002). In spherical coordinates $`(r,\theta ,\phi )`$, the $`(x,y,z)`$ components of the halo field are given by:
$`B_x=3\mu _\mathrm{G}\mathrm{sin}\theta \mathrm{cos}\theta \mathrm{cos}\phi /r^3`$
$`B_y=3\mu _\mathrm{G}\mathrm{sin}\theta \mathrm{cos}\theta \mathrm{sin}\phi /r^3`$ (10)
$`B_z=\mu _\mathrm{G}(13\mathrm{cos}^2\theta )/r^3`$
where $`\mu _\mathrm{G}184.2\mu \mathrm{G}\mathrm{kpc}^3`$ is the magnetic moment of the Galactic dipole. The dipole field is very strong in the central region of the Galaxy, but is only 0.3 $`\mu `$G in the vicinity of the Solar system, directed toward the North Galactic Pole.
There is a significant turbulent component, $`B_{\mathrm{random}}`$, of the Galactic magnetic field. Its field strength is difficult to measure but results found in literature are in the range of $`B_{\mathrm{random}}=0.5\mathrm{}2B_{\mathrm{reg}}`$ (Beck, 2001). However, we neglect the random field through the paper. Possible dependence of the results on the random field is discussed in the section 5.
## 4 NUMERICAL METHOD
### 4.1 Method of Calculation for Propagation of UHECRs
We numerically calculate an inverse process of propagation of UHE protons arriving at the earth in the intergalactic space. This method explained below is an expansion of many previous works on their propagation in the Galactic space (Fl$`\ddot{\mathrm{u}}`$ckiger et al., 1991; Bieber, Evenson, & Lin, 1992; Stanev, 1997; Medina-Tanco, 1999; Yoshiguchi et al., 2003b).
We already performed numerical simulations for UHECR propagation in the GMF in Yoshiguchi et al. (2003b). In the paper, we injected UHECRs from the earth isotropically, and recorded these trajectories until they reached a sphere of radius $`40`$ kpc centered at the Galactic center. The charge of UHECRs was taken as $`1`$ because we followed propagation of UHE protons backward. These UHECRs were injected with spectral index of $`2.7`$, which was similar to the observed one. Note that this is not the energy spectrum injected at extragalactic sources.
In this study, we expand these trajectories to the extragalactic space. In other words, there are our initial positions of UHECRs on a sphere of radius $`40`$ kpc centered at the Galactic center, which are the result of our previous work. The trajectories are followed until their distance from the Galaxy reaches 1 Gpc or their time for propagation reaches the age of the universe or their energies reach $`10^{25}`$ eV. Of course, we set the charge of UHECRs to be $`1`$.
In the extragalactic space, we have to consider not only the deflections due to the EGMF but also the energy loss processes (Berezinsky & Grigorieva, 1988; Yoshida & Teshima, 1993). UHE protons below $`4\times 10^{19}`$ eV lose their energies mainly due to adiabatic energy losses and pair production in collision with the cosmic microwave background (CMB). At the higher energies the photopion production with the CMB becomes essential (Detail explanation is given below). Though we assume that UHECRs are protons in this work, we should also add the photo-disintegration if we assume UHECRs to be nuclei. We treat all these energy loss processes as continuous processes. Note that energies of UHECRs increase during propagation because we follow their inverse processes.
The adiabatic energy loss is the effect of the expanding universe. This energy loss is written as
$$\frac{dE}{dt}=\frac{\dot{a}}{a}E=H_0\left[\mathrm{\Omega }_m(1+z)^3+\mathrm{\Omega }_\mathrm{\Lambda }\right]E.$$
(11)
As mentioned in the section 2, the cosmological parameters used in this calculation are $`\mathrm{\Omega }_m=0.3`$, $`\mathrm{\Omega }_\lambda =0.7`$, and $`H_0=71\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$.
The pair production due to collisions with the CMB can be treated as a continuous process which has small inelasticity ($`10^3`$). We adopt the analytical fit functions given by Chodorowski, Zdziarske, & Sikora (1992) to calculate the energy loss rate on isotropic photons.
UHE protons above $`4\times 10^{19}`$ eV lose a large fraction of their energy ($`20\%`$ at every reaction) in the photopion production with the CMB. We treated this process as a stochastic process in previous work. But in this study, we cannot treat this as a stochastic process because we calculate the inverse process. Berezinsky, Gazizov & Grigorieva (2002) and Aloisio & Berezinsky (2004) showed that the energy spectrum which is calculated with a continuous process of the photopion production is consistent with one calculated with a stochastic process if a distance between the earth and sources of UHECR is more than $`30`$ Mpc. Thus, we can adopt a continuous energy loss process since our source model (explained below) almost satisfy that condition about source distance. We use the energy loss length which is calculated by simulating the photopion production with the event generater SOPHIA (Mucke et al., 2000).
### 4.2 Source Distribution
We construct source models of UHECRs from our sample of galaxies explained in the section 3. The number density of UHECR sources is taken as our model parameter. For a given number density, we randomly select galaxies from our sample with probability proportional to absolute luminosity of each galaxy. We then estimate the source number density which reproduces the observed arrival distribution of UHECRs, by calculating the harmonic amplitude and the two point correlation function of arrival distribution of UHECRs as a function of the source number density.
### 4.3 Calculation of the UHECR Arrival Distribution
In this subsection, we explain the method of construction of UHECR arrival distribution at the earth.
We calculate 500,000 trajectories of UHE protons in the EGMF, using our method explained in Section 4.1, and record them. With our source models, we calculate a factor for each trajectory, which represents a relative probability that $`j`$ th proton reaches the Earth,
$$P_{\mathrm{selec}}(E,j)\underset{i}{}\frac{L_{i,j}}{(1+z_{ij})d_{i,j}^{}{}_{}{}^{2}}\frac{dN/dE_g(d_{i,j},E_i)}{E^{2.7}}\frac{dE_g}{dE}.$$
(12)
Here, $`i`$ labels sources on each trajectory. $`z_{i,j}`$, $`d_{i,j}`$ and $`L_{i,j}`$ is a redshift, a distance and luminosity of each source which is passed by $`j`$ th proton respectively. $`dN/dE(d_{i,j},E_i)`$ is the energy spectrum of protons at a source of distance $`d_{i,j}`$. $`E_i`$ is energy of proton at $`i`$ th source. $`E_g=E_g(E,d)`$ is the energy of cosmic ray at a source, which has the energy $`E`$ at the earth. $`dE_g/dE`$ represents variation of shape of the energy spectrum through propagation.
$`dE_g/dE`$ can be calculated in the case of rectilinear propagation (Berezinsky & Grigorieva, 1988). But calculation of this is difficult in this study since protons injected from the sources which are located at the same distance have different path length due to the EGMF. That is, the only $`E_g`$ cannot be decided when $`E`$ and $`d`$ are given. We calculate this factor using our 500,000 trajectories of protons.
Figure 4.3 shows a variation of shape of a monoenegetic spectrum ($`E=10^{19.6}`$ eV) at the earth for an example. The solid histogram is the spectrum at the earth. The dashed histogram and the dotted histogram are the spectra of UHECRs injected from the earth at 300 Mpc from us and 500 Mpc respectively. It is difficult to determine $`dE_g`$ from the figure since the spectra at far distances from the earth have large variances due to difference of their path lengthes. In this case, we calculate $`dE_g/dE`$ as
$$\frac{dE_g}{dE}(E,d)=\frac{dN/dE(E)}{dN/dE_g(E_{g}^{}{}_{}{}^{}(E,d),d)},$$
(13)
where $`dN/dE(E)`$ is the spectrum at the earth ($`E^{2.7}`$), $`dN/dE_g(E_g(E,d),d)`$ is a spectrum of UHECRs injected from the earth at a distance $`d`$ and $`E_{g}^{}{}_{}{}^{}(E,d)`$ is averaged $`E_g`$ when $`E`$ and $`d`$ are given.
We randomly select several trajectories according to these relative probabilities, so that the number of the selected trajectories is equal to the observed event number. The mapping of the velocity directions of each UHECR at the earth becomes the arrival distribution of UHECRs. The validity of this method is supported by the Liouville’s theorem.
If we have to select the same trajectory more than once in order to adjust the number of the selected trajectories, we generate a new event whose arrival angle is calculated by adding a normally distributed deviate with zero mean and variance equal to the experimental resolution $`2.8^{}`$ $`(1.8^{})`$ for $`E>10^{19}`$ eV $`(4\times 10^{19}\mathrm{eV})`$ to the original arrival angle.
We assume that UHECRs are protons injected with a power law spectrum in the range of $`10^{19}10^{22}`$ eV. We set this power law index 2.6 in order to fit the calculated energy spectrum to the one observed by the AGASA (Marco, Blasi, & Olinto, 2003). In other words, $`dN/dE(d_{i,j},E)E^{2.6}`$ where $`E`$ is the energy of UHECR at the source.
### 4.4 Statistical Methods
In this subsection, we explain the two statistical quantities, the harmonic analysis for large scale anisotropy (Hayashida et al., 1999) and the two point correlation function for small scale anisotropy.
The harmonic analysis to the right ascension distribution of events is the conventional method to search for large scale anisotropy of cosmic ray arrival distribution. For a ground-based detector like the AGASA, the almost uniform observation in right ascension is expected. The $`m`$-th harmonic amplitude $`r`$ is determined by fitting the distribution of cosmic rays to a sine wave with period $`2\pi /m`$. For a sample of $`n`$ measurements of phase, $`\varphi _1`$, $`\varphi _2`$, $`\mathrm{}`$, $`\varphi _n`$ (0 $`\varphi _i2\pi `$), it is expressed as
$$r=(a^2+b^2)^{1/2}$$
(14)
where, $`a=\frac{2}{n}\mathrm{\Sigma }_{i=1}^n\mathrm{cos}m\varphi _i`$, $`b=\frac{2}{n}\mathrm{\Sigma }_{i=1}^n\mathrm{sin}m\varphi _i`$. We calculate the harmonic amplitude for $`m=1,2`$ from a set of events generated by the method explained in the section 4.3.
If events with total number $`n`$ are uniformly distributed in right ascension, the chance probability of observing the amplitude $`r`$ is given by,
$$P=\mathrm{exp}(k),$$
(15)
where
$$k=nr^2/4.$$
(16)
The current AGASA 775 events above $`10^{19}`$ eV is consistent with isotropic source distribution within 90 $`\%`$ confidence level (Takeda et al., 2001). We therefore compare the harmonic amplitude for $`P=0.1`$ with the model prediction.
The two point correlation function $`N(\theta )`$ contains information on the small scale anisotropy. We start from a set of events generated from our simulation. For each event, we divide the sphere into concentric bins of angular size $`\mathrm{\Delta }\theta `$, and count the number of events falling into each bin. We then divide it by the solid angle of the corresponding bin, that is,
$`N(\theta )={\displaystyle \frac{1}{2\pi |\mathrm{cos}\theta \mathrm{cos}(\theta +\mathrm{\Delta }\theta )|}}{\displaystyle \underset{\theta \varphi \theta +\mathrm{\Delta }\theta }{}}1[\mathrm{sr}^1],`$ (17)
where $`\varphi `$ denotes the separation angle of the two events. $`\mathrm{\Delta }\theta `$ is taken to be $`1^{}`$ in this analysis. The AGASA data shows correlation at small angle $`(3^{})`$ with 2.3 (4.6) $`\sigma `$ significance of deviation from an isotropic distribution for $`E>10^{19}`$ eV $`(E>4\times 10^{19}\mathrm{eV})`$ (Takeda et al., 2001).
## 5 RESULTS
In this section, we present results of our simulations. In section 5.1, we constrain number density of UHECR sources from the observational results of the AGASA. Using our source model with this number density, we see how the EGMF affects the arrival distribution of UHECRs in section 5.2 and section 5.3.
### 5.1 A constraint on source model of UHECRs
In this subsection, we constrain source number density of UHECRs from the arrival distribution obtained by the AGASA.
Figure 4 and figure 5 show simulated harmonic amplitudes. The number of simulated events is set to be 775 in the energy above $`10^{19}`$ eV and their arrival direction is restricted in the range of $`10^{}\delta 80^{}`$ in order to compare our results with those of the AGASA. Note that $`\delta `$ is the declination. The shaded regions represent 1 $`\sigma `$ total statistical error, which is caused by the two components of statistical error which occur from the finite number of simulated events and the random source selection from our IRAS sample. In order to see magnitudes of each error, we also draw errorbars, which represent the only statistical fluctuation due to the finite number of the simulated events. The event selection and the random source selection are performed 100 times and 40 times respectively. The regions below the solid lines are expected for the statistical fluctuation of isotropic source distribution with the chance probability larger than 10%. For all source number density, both first and second amplitudes show that our source models predict sufficient isotropy of UHECR arrival distribution obtained by the AGASA.
Next, we investigate what number density of the sources reproduce the two point correlations obtained by the AGASA best. In order to evaluate it, we introduce $`\chi _{\theta _{\mathrm{max}}}`$ for a source distribution as
$$\chi _{\theta _{\mathrm{max}}}=\frac{1}{\theta _{\mathrm{max}}}\sqrt{\underset{\theta =0}{\overset{\theta _{\mathrm{max}}}{}}\frac{\left\{N(\theta )N_{\mathrm{obs}}(\theta )\right\}^2}{\sigma (\theta )^2}},$$
(18)
where $`N(\theta )`$ is the two point correlation function calculated from simulated arrival distribution within $`10^{}\delta 80^{}`$ and $`N_{\mathrm{obs}}(\theta )`$ is that obtained from the AGASA data at angle $`\theta `$. $`\sigma (\theta )`$ is total statistical error of $`N(\theta )`$ due to the finite number of simulated events. The random event selection are performed 100 times. This $`\chi _{\theta _{\mathrm{max}}}`$ represents goodness of fitting between the simulated two point correlation and the observed one. In this study, we take $`\theta _{\mathrm{max}}`$ to be 10.
Figure 5.1 shows $`\chi _{10}`$ as a function of number density of the sources ($`n_{\mathrm{source}}`$). The errorbars represent the statistical fluctuations due to the random source selection from our galaxy sample. The random source selections are performed 40 times. The two point correlation functions of the left and right panel are calculated in the energy range of $`E>4\times 10^{19}`$ eV and $`E>10^{19}`$ eV, respectively. However, we cannot know the normalization of two point correlation function obtained from the AGASA for $`E>10^{19}`$ eV. Thus, we fit the two point correlation function obtained by the AGASA data to that calculated from our simulation at $`\theta _{\mathrm{max}}`$.
Source models with larger source number density have strong peak at a small angle scale on the two point correlation since there are some sources near to the earth. On the other hand, in the case of smaller source number density, a small number of sources near to the earth contribute the arrival distribution, especially in the highest energy case (left panel), since radial distances between any two sources from the earth are more distant. Thus the peak of the two point correlation function also becomes more strong. Therefore $`\chi _{10}`$s should have a minimum as a function of source number density. As is seen from the figure, $`n_{\mathrm{source}}5\times 10^6`$ Mpc<sup>-3</sup> reproduces the two point correlation function obtained by the AGASA best.
Therefore, $`n_{\mathrm{source}}5\times 10^6`$ Mpc<sup>-3</sup> is the most appropriate number density of UHECR sources, since this source model also reproduces the harmonic amplitude obtained by the AGASA well. Note that number density of the sources have some uncertainty since the error bars in both panels are large.
Figure 5.1 is the energy spectra at the earth predicted by our source models. These spectra are averaged ones among 40 source distributions on each source number density. Solid line represents energy spectrum obtained for $`n_{\mathrm{source}}5\times 10^6`$ Mpc<sup>-3</sup>. We also show the observed cosmic-ray spectrum by the AGASA (Hayashida et al., 2000). These simulated energy spectra have cutoffs around $`E10^{19.68}`$ eV except $`n_{\mathrm{source}}1\times 10^5`$ Mpc<sup>-3</sup>, which are the GZK cutoff. These spectra can reproduce the AGASA data below $`10^{20}`$ eV. Note that the spectrum with $`n_{\mathrm{source}}1\times 10^5`$ Mpc<sup>-3</sup> has little spectral cutoff since there are many distributions which contain sources in the GZK sphere due to large source number density. The source model ($`n_{\mathrm{source}}5\times 10^6`$ Mpc<sup>-3</sup>) also reproduces the observed energy spectrum only below $`10^{20}`$eV. We concluded in Yoshiguchi et al. (2003a) that a large fraction of cosmic rays above $`10^{20}`$eV observed by the AGASA might originate in the top-down scenarios. Thus we consider UHECRs with only $`E<10^{20}`$eV in what follows.
### 5.2 Arrival Distribution of UHECRs above $`10^{19}`$eV
In this subsection, we demonstrate a skymap of the arrival distribution of UHECRs in the case of $`n_{\mathrm{source}}5\times 10^6`$ Mpc<sup>-3</sup>. We construct their arrival distribution using our method explained in the section 4.3.
Figure 6 shows one of results of the event generation above $`10^{19}`$eV calculated from a specific source model with $`n_{\mathrm{source}}5\times 10^6`$ Mpc<sup>-3</sup>. The points represent each event and the events are colored according to their energies. The number of events is 5000, which is the expected number of events observed by the Pierre Auger observatory for a few years (Capelle et al., 1998). The skymap generated with both the EGMF and the GMF is in the lower right panel and that without any magnetic fields, that with only the EGMF and that with only the GMF is in the upper left, the upper right and the lower left panel respectively.
This specific source model has three strong sources (see upper right). One is $`(l,b)(199^{},34^{})`$, another is $`(l,b)(287^{},19^{})`$ and the other is $`(l,b)(25^{},11^{})`$. Each distance from us is about 77 Mpc, 65 Mpc and 70 Mpc respectively. In the absence of any magnetic fields (upper left panel), there are the strong clusterings of events at the directions of these three sources. When the effects of the EGMF are included (upper right), we find the diffusion of the clustered events. In the lower left panel, the clustered events are arranged in the order of their energies, reflecting the directions of the GMF. This was pointed out by Alvarez-Muniz, Engel & Stanev (2002) and Yoshiguchi et al. (2003b). Note that we cannot find the clustered events at direction of one of the strong sources $`(l,b)=(25^{},11^{})`$. This is because UHE protons injected at this source cannot reach the earth due to the GMF. In the lower right panel, we also find the arrangements at the same points of the lower left panel. In addition, the EGMF diffuses these clustered events as we see in the two upper panels.
In order to see these features quantitatively, we compare the statistical quantities calculated with the EGMF to those calculated without the EGMF in the presence of the GMF in the next subsection.
### 5.3 Statistics on the UHECR Arrival Distribution
In this subsection, we compare the statistical quantities on the arrival distribution calculated with the EGMF to those without the EGMF. We take $`n_{\mathrm{source}}5\times 10^6`$ Mpc<sup>-3</sup>.
Figure 5.3 and figure 5.3 show the two point correlation functions simulated by our source model in the energy range of $`E>4\times 10^{19}`$eV and $`E>10^{19}`$eV respectively. In each figure, the left panel shows the two point correlation function calculated with the EGMF and the right panel shows that without the EGMF. Note that the GMF is taken into account in the figures. We calculate the two point correlation function for the simulated events within only $`10^{}\delta 80^{}`$ in order to compare our results with the AGASA data. The shaded regions represent 1 $`\sigma `$ total statistical error, which is caused by the finite number of simulated events and the random source selections from our IRAS sample. We also draw errorbars, which represent the statistical fluctuations due to the finite number of the simulated events, which is set to be equal to that observed by the AGASA (49 events for $`E>4\times 10^{19}`$ eV and 775 events for $`E>10^{19}`$ eV). The event selection and the source selection are performed 100 times and 40 times respectively. The histograms represent the AGASA data (Takeda et al., 2001). For the data of $`E>10^{19}`$eV, we normalize the two point correlation function as the correlation function obtained by the AGASA fits the calculated one at 30, since we cannot know the normalization of the AGASA data with this energy.
In both figure 5.3 and 5.3, it is visible that a peak at small angle is much stronger than that of the AGASA though the AGASA data were covered in the shaded regions, which are mainly caused by the source selection. In our previous work (Yoshiguchi et al., 2003b) in which the EGMF was neglected, we also faced this situation and pointed out possible explanations, one of which was effects of the EGMF.
In figure 5.3, we find that a peak of two point correlation function calculated with the EGMF at small angle scale is weaker than that without the EGMF. Thus the consistency with the AGASA data becomes better due to the EGMF. On the other hand, in figure 5.3, a obvious difference between the two panels on the peak at small angle is not found. This is because UHECRs above $`4\times 10^{19}`$ eV are less deflected by the EGMF and hardly dispersed.
As mentioned above, the calculated two point correlation functions have large errors since some source distributions out of 40 contain very near source in our source model. Such source distributions do not reproduce the large-scale isotropy observed by the AGASA. We check that 20 source distributions out of 40 predict the sufficient large-scale isotropy obtained by the AGASA within 1 $`\sigma `$ statistical error due to the finite number of the simulated events. From these 20 source distributions, we calculate the two point correlation function in the presence of the EGMF again. The results are shown in figure 5.3.
In figure 5.3, we find that the consistency with the AGASA data becomes better than the two point correlation functions in figure 5.3 and 5.3. We also see that the errors due to the source selection become small. However, we should note that the peak at small angle scale is still relatively strong, compared with the AGASA though our previous result (Yoshiguchi et al., 2003b) is improved by the effect of the EGMF. We assume effects of the random component of the GMF, which is neglected in this work, as one of possible explanations for this fact. This issue is left for future investigations.
We also investigate the harmonic amplitudes in the same way. But there is little difference dependent on the EGMF.
## 6 SUMMARY AND DISCUSSION
We presented numerical simulations on propagation of UHECRs above $`10^{19}`$ eV in a structured EGMF and GMF. We used the IRAS PSCz catalogue in order to construct a structured EGMF model which reflects the local structures actually observed, and source models of UHECRs. We calculated an inverse process of their propagation taking the energy loss processes into account in the EGMF. We injected UHECRs from the earth isotropically whose charges are taken as -1 and recorded these trajectories. They could be regarded as trajectories of UHE protons from the extragalactic space. We then select some of their trajectories according to given source distributions. The simulated arrival distribution was able to be obtained by mapping the velocity directions of the selected trajectories at the earth. The use of this method enabled us to calculate only trajectories of UHECRs reaching the earth and saved the CPU time effectively. The validity of this method was supported by the Liouville’s theorem.
We calculated the harmonic amplitudes and the two point correlation functions of arrival distribution of UHECRs above $`10^{19}`$eV, using our source models and examined what number density of the sources reproduces the large-scale isotropy and the small-scale anisotropy obtained by the AGASA best. As a result, we found that $`5\times 10^6`$ Mpc<sup>-3</sup> was the most appropriate number density of source of UHECRs. Number density of source is a constraint on source candidate of UHECRs.
We also demonstrated skymaps of the arrival distribution of UHECRs above $`10^{19}`$ eV, using our source model for $`n_{\mathrm{source}}5\times 10^6`$ Mpc<sup>-3</sup> with the event number expected by future experiments and examined how the EGMF affects the arrival distribution of UHECRs. The main result was diffusion of clustering events which is obtained by calculations neglecting the EGMF.
In order to see the effect of the EGMF quantitatively, we compared the two point correlation function calculated with our structured EGMF model to that without the EGMF. We found that the EGMF weakened the small-scale anisotropy and improved a prediction in Yoshiguchi et al. (2003b), which had been calculated with only the GMF.
However, we found the calculated two point correlation functions had large errorbars since source distributions, which contained sources very near to us, existed. Such source distributions do not reproduce the large-scale isotropy observed by the AGASA. Thus we calculated the two point correlation functions from source distributions which predicted the large-scale isotropy obtained by the AGASA within 1 $`\sigma `$ statistical error. These simulated two point correlation function reproduced that of the AGASA better. It is possible that these functions at small angle scale can be closer to the observational data due to the random component of the GMF. This issue is left for future studies.
New large aperture detectors are under development, such as the Pierre Auger observatory (Capelle et al., 1998), the EUSO (Benson & Linsley, 1982) and the Telescope Array. These projects are expected to increase observed events of UHECRs per year drastically. We can obtain more strong constraints on our source model, other than number density of sources, using other statistical quantities when the detailed data of large events of UHECRs are published. This is one of plans of future studies.
The work of H.Y. is supported by Grants-in-Aid for JSPS Fellows. The work of K.S. is supported by Grants-in-Aid for Scientific Research provided by the Ministry of Education, Science and Culture of Japan through Research Grant No.S14102004 and No.14079202.
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# Matrix Factorizations and Representations of Quivers I
## 1. Introduction
This paper introduces new triangulated categories associated to quasi-homogeneous polynomials which define isolated singularities only at the origin and relates those categories with the derived categories of representations of quivers. Our motivation comes from K. Saito’s theory of primitive forms, especially from a problem in his study on regular weight systems and generalized root systems \[Sa1\]. We will explain the problem below.
Let $`(a,b,c;h)`$ be a quadruple of positive integers such that the function
$$\chi (T):=\frac{(T^{ha}1)(T^{hb}1)(T^{hc}1)}{(T^a1)(T^b1)(T^c1)}$$
has no poles. Such a quadruple $`W:=(a,b,c;h)`$ is called a regular weight system. It is known that $`W`$ is a regular weight system if and only if we have at least one polynomial $`f(x,y,z)[x,y,z]`$ such that
$$ax\frac{f}{x}+by\frac{f}{y}+cz\frac{f}{z}=hf$$
and
$$X_0:=\{(x,y,z)^3|f(x,y,z)=0\}$$
has an isolated singularity only at the origin. Note that the (restricted) map
$$f:^3\backslash f^1(0)\backslash \{0\}$$
is a topologically locally trivial fiber bundle, and the general fiber $`X_1:=f^1(1)`$ (called the Milnor fiber) is an open 2-dimensional complex manifold whose second homology group $`H_2(X_1,)`$ is a free $``$-module of rank $`\mu :=(ha)(hb)(hc)/abc=lim_{t1}\chi (T)`$. Since $`X_1`$ is real 4-dimensional, $`H_2(X_1,)`$ has an intersection form
$$I_{H_2(X_1,)}:H_2(X_1,)\times H_2(X_1,).$$
If $`f`$ is a defining polynomial of a simple (ADE) singularity, then $`(H_2(X_1,),I_{H_2(X_1,)})`$ gives the root lattice of the finite root system corresponding to the singularity. In \[Sa1\], it is shown that $`(H_2(X_1,),I_{H_2(X_1,)})`$ with the set of vanishing cycles (which corresponds to the set of roots) and the Milnor monodromy (which corresponds to the Coxeter transformation) satisfies the axioms of the generalized root system which naturally extends the classical (finite) root systems. Since both weight systems and generalized root systems are combinatorial, it is natural to propose the following problem.
###### Problem 1.1.
$`(`$\[Sa1\]$`)`$
Construct directly from a regular weight system $`W`$, without passing through the homology group $`H_2(X_1,)`$ of the Milnor fiber, arithmetically or combinatorially, the generalized root system of the vanishing cycles.
The purpose of this paper is to develop a necessary tools in terms of $`A_{\mathrm{}}`$-categories and to give a partial answer to the above problem. Let $`k`$ be a field of characteristic zero. First, we introduce a notion of $``$-graded $`A_{\mathrm{}}`$-categories over $`k`$ (Definition 2.1) in order to consider the polynomial ring $`[x_1,\mathrm{}x_n]`$ with a polynomial $`f`$ satisfying the quasi-homogeneous condition
$$\underset{i=1}{\overset{n}{}}\frac{2a_i}{h}x_i\frac{f}{x_i}=2f,$$
(1.1)
where $`a_1,\mathrm{},a_n`$ and $`h`$ are positive integers such that the greatest common divisor of them is $`1`$, as a usual $`_2`$-graded $`A_{\mathrm{}}`$-category with $`m_0(1)=f`$ and an ”extra $`\frac{2}{h}`$-grading”. We shall denote the $``$-graded $`A_{\mathrm{}}`$-category defined by $`f[x_1,\mathrm{},x_n]`$ by $`𝒜_f`$ (Example-Definition 2.3).
Next, we consider the category of twisted complexes over $``$-graded $`A_{\mathrm{}}`$-categories (Proposition 2.15) and the derived category of $``$-graded $`A_{\mathrm{}}`$-categories (Definition 2.17). The important fact is that the twisted complexes over $`𝒜_f`$ coincide with matrix factorizations of $`f`$ introduced by Eisenbud \[E\] in his study of maximal Cohen-Macaulay modules. Since we consider the quasi-homogeneous polynomial $`f`$, we have a group action ($``$-action) on the category of matrix factorizations. Inspired by the work by Hori and Walcher \[HW\], we introduce the $``$-equivariant derived category of $`𝒜_f`$ denoted by $`D_{}^b(𝒜_f)`$ (Definition 2.23). We can now propose a conjecture to K. Saito’s problem.
###### Conjecture 1.2.
Let $`W`$ be a regular weight system and $`f`$ be a quasi-homogeneous polynomial attached to $`W`$. Assume $`W`$ has a dual regular weight system $`W^{}=(a^{},b^{},c^{};h)`$ in the sense of \[Sa2\] and let $`f^{}`$ be a quasi-homogeneous polynomial attached to $`W^{}`$. Then the following should hold.
1. $`D_{}^b(𝒜_f^{})`$ is generated as a triangulated category by objects $`\{E_1,\mathrm{},E_\mu \}`$ such that
$$\mathrm{Hom}_{D_{}^b(𝒜_f^{})}(E_i,E_j)=0,\text{if }i>j,\mathrm{Hom}_{D_{}^b(𝒜_f^{})}(E_i,E_j[k])=0,k0,{}_{}{}^{}i,j.$$
(1.2)
That is to say, $`D_{}^b(𝒜_f^{})`$ is generated by a strongly exceptional collection.
2. $`D_{}^b(𝒜_f^{})`$ has the Serre functor $`S`$ such that $`S^h[3h2a^{}2b^{}2c^{}]`$ where $`[1]`$ is the shift functor on $`D_{}^b(𝒜_f^{})`$.
3. Let $`a_{ij}:=\chi (E_i,E_j)=dim_k\mathrm{Hom}_{D_{}^b(𝒜_f^{})}(E_i,E_j)`$. Put $`A:=(a_{ij})`$ and $`I_{K_0(D_{}^b(𝒜_f^{}))}=A^1+{}_{}{}^{t}A_{}^{1}`$. Then $`(K_0(D_{}^b(𝒜_f^{})),I_{K_0(D_{}^b(𝒜_f^{}))})`$ is isomorphic to $`(H_2(X_1,),I_{H_2(X_1,)})`$ as a lattice.
This conjecture is based on the relation between the duality of regular weight systems and the mirror symmetry of Landau-Ginzburg orbifolds (see \[T\]). We do not discuss this background in detail here but we write the following diagram for reader’s convenience.
$$\begin{array}{ccc}\text{Quasi-homogeneous polynomial }f\text{ for }W& \stackrel{Milnorfiber}{}& \{\text{Vanishing cycles in }X_1=f^1(1)\}\\ \text{Duality of weights + Orbifold}& & ||\\ \{\text{B-branes in LG orbifold }W^{}//(/h)\}& \stackrel{MirrorSymmetry}{}& \{\text{A-branes in LG model for }W\}\end{array}$$
For ADE singularities, we know that $`WW^{}`$ and the generalized root systems for them are the classical finite root systems. Therefore, we may expect that the following conjecture should hold.
###### Conjecture 1.3.
Let $`W`$ be a regular weight system corresponding to an ADE singularity and $`f`$ be a quasi-homogeneous polynomial attached to $`W`$. Then $`D_{}^b(𝒜_f)`$ is equivalent as a triangulated category to the bounded derived category of finite dimensional representations of Dynkin quiver corresponding to the type of singularity of $`f`$.
This is also expected from the homological mirror symmetry phenomena for ADE singularities studied by Seidel. See \[Se2\] for details.
In this paper, we will prove the conjecture for $`A_n`$-singularities (Theorem 3.1), where we reduce to the case $`f:=x^{n+1}[x]`$ by Knörrer’s periodicity \[K\] (see also \[O1\]). We will give a proof of the above conjecture for general cases in a separate paper \[KST\].
Finally, we will construct a special stability condition for the triangulated category $`D_{}^b(𝒜_f)`$ in the sense of T. Bridgeland \[B\] for $`f=x^h`$. We can naturally introduce in our formulation the phase of objects (Definition 4.1) and the central charge $`Z_\omega `$ (Definition 4.3).
While our preparation of this paper, two papers related to our work appeared. One is the paper \[W\] by J. Walcher where he studies from physical point of view the similar categories and the stability conditions on them (his notion of ”R-stability”). Another is the paper \[O2\] by D. Orlov where he studies the triangulated category for singularities with a $`^{}`$-action, which is equivalent to our category $`D_{}^b(𝒜_f)`$.
Acknowledgement
I am grateful to Kentaro Hori, Hiroshige Kajiura, Masaki Kashiwara and Kyoji Saito for very useful discussions. This work was partly supported by Grant-in Aid for Scientific Research grant numbers 14740042 of the Ministry of Education, Science and Culture in Japan.
## 2. $``$-graded $`A_{\mathrm{}}`$-categories
In this section, we set up several definitions which we will use in the later sections. Let $`k`$ be a field of characteristic zero.
###### Definition 2.1.
Let $`h`$ be a positive number. A $``$-graded $`A_{\mathrm{}}`$-category $`𝒜`$ of index $`h`$ is a collection of the following data.
1. A set of objects $`Ob(𝒜)`$,
2. A set of homomorphisms, a $`\times _2`$ graded $`k`$-linear vector space for each $`a,bOb(𝒜)`$
$$𝒜(a,b)=\underset{q}{}𝒜^q(a,b)_+𝒜^q(a,b)_{},$$
such that
$$𝒜^q(a,b)_+=0,q\frac{2}{h},𝒜^q(a,b)_{}=0,q1\frac{2}{h}.$$
We call the subspaces
$$𝒜(a,b)_+:=\underset{q\frac{2}{h}}{}𝒜^q(a,b)_+,𝒜(a,b)_{}:=\underset{q1\frac{2}{h}}{}𝒜^q(a,b)_{}$$
the even and the odd subspaces.
3. for $`n0`$, $`k`$-multilinear maps
$$m_n^𝒜:𝒜(a_{n1},a_n)\mathrm{}𝒜(a_0,a_1)𝒜(a_0,a_n),a_iOb(𝒜),$$
of degree $`2n`$ with respect to the $``$-grading which is even $`(`$odd$`)`$ with respect to the $`_2`$-grading when $`n`$ is even $`(`$odd$`)`$, where $`m_0^𝒜`$ is a map
$$m_0^𝒜:k𝒜(a,a).$$
The multilinear maps satisfy the following $`(A_{\mathrm{}}`$-relation$`)`$. For fixed $`n`$, we have
$$\begin{array}{c}\underset{r+s+t=n}{}\underset{r+1+t=u}{}(1)^{|x_1|+\mathrm{}+|x_r|+r}m_u^𝒜(x_{r+s+t}\mathrm{}x_{r+s+1}\hfill \\ \hfill m_s^𝒜(x_{r+s}\mathrm{}x_{r+1})x_r\mathrm{}x_1)=0,\end{array}$$
(2.1)
where $`|x_i|`$ is the parity of the morphism defined by
$$|x_i|:=\{\begin{array}{c}0,x_i𝒜(a_{i1},a_i)_+,\hfill \\ 1,x_i𝒜(a_{i1},a_i)_{}\hfill \end{array}.$$
(2.2)
###### Remark 2.2.
$``$-graded $`A_{\mathrm{}}`$-category of index $`1`$ is nothing but an $`A_{\mathrm{}}`$-category with the usual $``$-grading, we call it a $``$-graded $`𝒜_{\mathrm{}}`$-category or simply an $`𝒜_{\mathrm{}}`$-category. See \[F\] and \[Se1\] for details of homological algebra of $`𝒜_{\mathrm{}}`$-categories.
We write down explicitly the relation (2.1) when $`m_n^𝒜=0`$ for $`n3`$. For $`x,y,z_{a,b}𝒜(a,b)`$, we have
$$m_1^𝒜(m_0^𝒜(1))=0,$$
$$m_1^𝒜(m_1^𝒜(x))=(1)^{|x|}m_2^𝒜(m_0^𝒜(1)x)m_2^𝒜(xm_0^𝒜(1)),$$
$$m_1^𝒜(m_2^𝒜(xy))=(1)^{|y|}m_2^𝒜(m_1^𝒜(x)y)m_2^𝒜(xm_1^𝒜(y)),$$
$$m_2^𝒜(m_2^𝒜(xy)z)=(1)^{|z|}m_2^𝒜(xm_2^𝒜(yz)).$$
Put $`u:=m_0^𝒜(1)`$, $`d(x):=(1)^{|x|+1}m_1^𝒜(x)`$ and $`xy:=(1)^{|y|}m_2^𝒜(xy)`$. Then a triple $`(u,d,)`$ defines on $`_{a,b}𝒜(a,b)`$ a curved differential graded (CDG) algebra structure \[KL1\].
###### Example-Definition 2.3.
Let $`f[x_1,\mathrm{}x_n]`$ be a polynomial which satisfies the following quasi-homogeneous condition$`:`$
$$\underset{i=1}{\overset{n}{}}\frac{2a_i}{h}x_i\frac{f}{x_i}=2f,$$
(2.3)
where $`a_1,\mathrm{},a_n`$ and $`h`$ are positive integers such that the greatest common divisor of them is $`1`$. We denote by $`𝒜_f`$ the $``$-graded $`A_{\mathrm{}}`$-category of index $`h`$ defined as follows:
$$Ob(𝒜_f)=\{a\},$$
$$𝒜_f(a,a):=[x_1,\mathrm{}x_n],𝒜_f(a,a)_{}:=0,$$
$$m_0^{𝒜_f}(1):=f𝒜_f^2(e,e)_+,m_1^{𝒜_f}:=0,$$
$$m_2^{𝒜_f}(\alpha \beta ):=\alpha \beta ,\alpha ,\beta [x_1,\mathrm{}x_n],$$
where $``$ is the usual product on $`[x_1,\mathrm{}x_n]`$.
In the above example, we have the special element $`1𝒜_f^0(a,a)_+`$ which defines a unit of the algebra $`[x_1,\mathrm{}x_n]`$. It is well-known that the notion of units in $`A_{\mathrm{}}`$-categories can be introduced as follows:
###### Definition 2.4.
Let $`aOb(𝒜)`$. $`e_a𝒜^0(a,a)_+`$ is called a unit if
$$m_2^𝒜(x,e_a)=x,m_2^𝒜(e_a,y)=(1)^{|y|}y,$$
(2.4)
for $`x𝒜(a,b)`$ and $`y𝒜(b,a)`$, and for $`n2`$
$$m_n^𝒜(x_1,\mathrm{},x_n)=0,$$
(2.5)
if one of $`x_i`$ coincides with $`e_a`$.
###### Remark 2.5.
It is easy to check that if a unit exists then it is unique.
###### Definition 2.6.
$``$-graded $`A_{\mathrm{}}`$-category is called unital if each object has a unit.
Let $`𝒜`$ be a $``$-graded $`A_{\mathrm{}}`$-category of index $`h`$. We can construct another $``$-graded $`A_{\mathrm{}}`$-category $`\overline{𝒜}`$ of index $`h`$ from $`𝒜`$ as follows:
###### Definition 2.7.
Let $`𝒜`$ be a unital $``$-graded $`A_{\mathrm{}}`$-category of index $`h`$.
1. a set of objects $`Ob(\overline{𝒜})`$ is given by
$$Ob(\overline{𝒜}):=\{a\{\frac{2k}{h}\},aOb(𝒜),k,b\{\frac{2l}{h}\}[1],bOb(𝒜),l\},$$
(2.6)
2. a set of homomorphisms is given by
$$\overline{𝒜}^q(a\{\frac{2k}{h}\},b\{\frac{2l}{h}\})_\pm :=𝒜^{q+\frac{2(lk)}{h}}(a,b)_\pm ,$$
(2.7)
$$\overline{𝒜}^q(a\{\frac{2k}{h}\},b\{\frac{2l}{h}\}[1])_\pm :=𝒜^{q+\frac{2(lk)}{h}1}(a,b)_{},$$
(2.8)
$$\overline{𝒜}^q(a\{\frac{2k}{h}\}[1],b\{\frac{2l}{h}\})_\pm :=𝒜^{q+\frac{2(lk)}{h}+1}(a,b)_{},$$
(2.9)
$$\overline{𝒜}^q(a\{\frac{2k}{h}\}[1],b\{\frac{2l}{h}\}[1])_\pm :=𝒜^{q+\frac{2(lk)}{h}}(a,b)_\pm .$$
(2.10)
3. $`k`$-multilinear maps $`m_n^{\overline{𝒜}}`$ are defined using those on $`𝒜`$ with additional signs as follows. For $`x_1\overline{𝒜}(a_0,a_1),\mathrm{},x_n\overline{𝒜}(a_{n1},a_n)`$,
$$m_n^{\overline{𝒜}}(x_n\mathrm{}x_1):=(1)^{|a_0|}m_n^𝒜(x_n\mathrm{}x_1),$$
where we regard $`x_i`$ in the right hand side as a homomorphism in $`𝒜`$ by the above definition $`(ii)`$ and $`|a_0|`$ is the parity of $`a_0`$ defined by
$$|a_0|:=\{\begin{array}{c}0,a_0=a\{\frac{2k}{h}\},aOb(𝒜),k\hfill \\ 1,a_0=a\{\frac{2l}{h}\}[1],aOb(𝒜),l\hfill \end{array}.$$
(2.11)
$`A_{\mathrm{}}`$-functors for $``$-graded $`A_{\mathrm{}}`$-categories can be defined in an obvious way. There are the ”translation” functor $`\{\frac{2}{h}\}`$ and the shift functor $`[1]`$ on $`\overline{𝒜}`$.
###### Proposition 2.8.
The following functors $`\{\frac{2}{h}\}`$ and $`[1]`$ define autoequivalences of $`\overline{𝒜}:`$
$$\{\frac{2}{h}\}\left(a\{\frac{2k}{h}\}\right):=a\{\frac{2(k+1)}{h}\},\{\frac{2}{h}\}\left(a\{\frac{2k}{h}\}[1]\right):=a\{\frac{2(k+1)}{h}\}[1],$$
(2.12)
and
$$[1]\left(a\{\frac{2k}{h}\}\right):=a\{\frac{2(k+h)}{h}\}[1],[1]\left(a\{\frac{2k}{h}\}[1]\right):=a\{\frac{2k}{h}\}.$$
(2.13)
Put $`\{\frac{2k}{h}\}:=\{\frac{2}{h}\}^k`$ and $`[l]:=[1]^l`$ for $`k,l`$. We have the relation $`\{\frac{2h}{h}\}=[2]`$. ∎
Let $`\overline{𝒜}`$ be an $``$-graded $`A_{\mathrm{}}`$-category of index $`h`$. Consider the $``$-graded $`A_{\mathrm{}}`$-category $`\stackrel{~}{𝒜}`$ of index $`h`$ whose set of objects is the set of finite (formal) direct sums of objects of $`\overline{𝒜}`$,
$$Ob(\stackrel{~}{𝒜}):=\left\{a=\underset{i}{}a_i\{\frac{2k_i}{h}\}\underset{j}{}a_j\{\frac{2l_j}{h}\}[1],a_i,a_jOb(\overline{𝒜}),k_i,l_j\right\},$$
(2.14)
whose set of homomorphisms is
$$\begin{array}{c}\stackrel{~}{𝒜}(a,b):=\underset{i_1,i_2}{}\overline{𝒜}(a_{i_1}\{\frac{2k_{i_1}}{h}\},b_{i_2}\{\frac{2k_{i_2}}{h}\})\underset{i_1,j_2}{}\overline{𝒜}(a_{i_1}\{\frac{2k_{i_1}}{h}\},b_{j_2}\{\frac{2l_{j_2}}{h}\}[1])\hfill \\ \hfill \underset{j_1,i_2}{}\overline{𝒜}(a_{j_1}\{\frac{2l_{j_1}}{h}\}[1],b_{i_2}\{\frac{2k_{i_2}}{h}\})\underset{j_1,j_2}{}\overline{𝒜}(a_{j_1}\{\frac{2l_{j_1}}{h}\},b_{j_2}\{\frac{2l_{j_2}}{h}\}[1]),\end{array}$$
(2.15)
and whose $`k`$-linear maps are defined by those on $`\overline{𝒜}`$ using the natural ”matrix multiplication” rule.
###### Definition 2.9.
Take an object $`aOb(\stackrel{~}{𝒜})`$ and $`Q\stackrel{~}{𝒜}(a,a)_{}`$. $`(a;Q)`$ is called a twisted complex if $`Q`$ satisfies the Maurer-Cartan equation
$$\underset{n0}{}m_n^{\stackrel{~}{𝒜}}(Q^n)=0.$$
(2.16)
The set of all twisted complexes is denoted by $`Ob(Tw(𝒜))`$. If $`Q\stackrel{~}{𝒜}^1(a,a)_{}`$ in addition, then $`(a;Q)`$ is called a graded twisted complex and we denote the set of all graded twisted complexes by $`Ob(Tw_{}(𝒜))`$.
An assumption is necessary for the equation (2.16) to make sense. If $`m_0^𝒜=0`$, then it is usually introduced that the notion of one-sided twisted complexes which makes the sum in the equation (2.16) finite. However we study in this paper the case when $`m_0^𝒜0`$, we shall assume that our $`A_{\mathrm{}}`$-categories have no higher product, in other words, $`m_n^𝒜=0`$ for all $`n3`$.
We often write $`Q`$ in the following form:
$$Q=\left(\begin{array}{cc}Q_{++}& Q_+\\ Q_+& Q_{},\end{array}\right)$$
(2.17)
where
$$Q_{\pm \pm }\stackrel{~}{𝒜}(a_\pm ,a_\pm )_{},Q_\pm \stackrel{~}{𝒜}(a_\pm ,a_{})_+,$$
(2.18)
and $`a_\pm `$ are given by the following decomposition
$$a=a_++a_{}[1],a_+=\underset{i}{}a_{+,i}\{\frac{2k_i}{h}\},a_{}=\underset{i}{}a_{,i}\{\frac{2l_i}{h}\}.$$
(2.19)
###### Remark 2.10.
If $`𝒜`$ is a unital $``$-graded $`A_{\mathrm{}}`$-category of index $`h`$ with $`m_n^𝒜=0`$, $`n3`$, then there exists at least one twisted complex for each object $`aOb(\stackrel{~}{𝒜})`$. Indeed,
$$Q|_{ij}=\{\begin{array}{c}0,\text{for }i,j\text{such that }a_{+,i}a_{,j}\text{ and }a_{+,i}\{\frac{2h}{h}\}a_{,j},\hfill \\ \left(\begin{array}{cc}0& m_0^𝒜(1)\\ e_{a_{+,i}}& 0\end{array}\right),\text{for }i,j\text{such that }a_{+,i}=a_{,j},\hfill \\ \left(\begin{array}{cc}0& e_{a_{+,i}}\\ m_0^𝒜(1)& 0\end{array}\right),\text{for }i,j\text{such that }a_{+,i}\{\frac{2h}{h}\}=a_{,j},\hfill \end{array}$$
(2.20)
is a twisted complex.
###### Example 2.11.
Since $`𝒜_f`$ has no odd homomorphisms, each twisted complex $`(a=a_+a_{}[1];Q_a)`$ has the following form
$$Q_a:=\left(\begin{array}{cc}0& Q_+\\ Q_+& 0\end{array}\right),Q_+\stackrel{~}{𝒜}(a_+,a_{})_+,Q_+\stackrel{~}{𝒜}(a_{},a_+)_+.$$
The Maurer-Cartan equation (2.16) becomes
$$f\mathrm{Id}Q_a^2=0.$$
(2.21)
This is exactly the same equation which first studied by Eisenbud \[E\] in his work on maximal Cohen-Macaulay modules. $`Q_a`$ is called a matrix factorization of $`f`$.
Let $`𝒜`$ be a unital $``$-graded $`A_{\mathrm{}}`$-category of index $`h`$ with $`m_n^𝒜=0`$, $`n3`$.
###### Definition 2.12.
Let $`\alpha :=(a;Q_a)`$ and $`\beta :=(b;Q_b)`$ be twisted complexes. We first put
$$Tw(𝒜)(\alpha ,\beta ):=\stackrel{~}{𝒜}(a,b)_+\stackrel{~}{𝒜}(a,b)_{}.$$
(2.22)
We define a $`k`$-multilinear maps $`m_n^{Tw(𝒜)}(n=0,1,2)`$ by
$$m_0^{Tw(\overline{𝒜})}(1):=0,$$
(2.23)
$$m_1^{Tw(𝒜)}(\mathrm{\Phi }):=m_1^{\stackrel{~}{𝒜}}(\mathrm{\Phi })+m_2^{\stackrel{~}{𝒜}}(Q_b\mathrm{\Phi })+m_2^{\stackrel{~}{𝒜}}(\mathrm{\Phi }Q_a),$$
(2.24)
where $`\mathrm{\Phi }Tw(𝒜)(\alpha ,\beta )`$ and
$$m_2^{Tw(𝒜)}(\mathrm{\Psi }_2\mathrm{\Psi }_1):=m_2^{\stackrel{~}{𝒜}}(\mathrm{\Psi }_2\mathrm{\Psi }_1),$$
(2.25)
for $`\mathrm{\Psi }_1Tw(𝒜)(\alpha _0,\alpha _1)=\stackrel{~}{𝒜}(a_0,a_1)`$ and $`\mathrm{\Psi }_2Tw(𝒜)(\alpha _1,\alpha _2)=\stackrel{~}{𝒜}(a_1,a_2)`$.
We often write the spaces of morphisms in the matrix form:
$$Tw(𝒜)(\alpha ,\beta )_\pm =\left(\begin{array}{cc}\stackrel{~}{𝒜}(a_+,b_+)_\pm & \stackrel{~}{𝒜}(a_{},b_+)_{}\\ \stackrel{~}{𝒜}(a_+,b_{})_{}& \stackrel{~}{𝒜}(a_{},b_{})_\pm \end{array}\right).$$
###### Lemma 2.13.
$`(m_1^{Tw(𝒜)})^2=0`$.
###### Proof.
For $`\mathrm{\Phi }_\pm Tw(𝒜)(\alpha ,\beta )_\pm `$, we have
$`(m_1^{Tw(𝒜)})^2(\mathrm{\Phi }_\pm )=`$ $`(m_1^{\stackrel{~}{𝒜}})^2(\mathrm{\Phi }_\pm )+m_1^{\stackrel{~}{𝒜}}(m_2^{\stackrel{~}{𝒜}}(Q_b\mathrm{\Phi }_\pm ))+m_1^{\stackrel{~}{𝒜}}(m_2^{\stackrel{~}{𝒜}}(\mathrm{\Phi }_\pm Q_a))`$
$`+m_2^{\stackrel{~}{𝒜}}(Q_b(m_1^{\stackrel{~}{𝒜}}(\mathrm{\Phi }_\pm )+m_2^{\stackrel{~}{𝒜}}(Q_b\mathrm{\Phi }_\pm )+m_2^{\stackrel{~}{𝒜}}(\mathrm{\Phi }_\pm Q_a)))`$
$`+m_2^{\stackrel{~}{𝒜}}((m_1^{\stackrel{~}{𝒜}}(\mathrm{\Phi }_\pm )+m_2^{\stackrel{~}{𝒜}}(Q_b\mathrm{\Phi }_\pm )+m_2^{\stackrel{~}{𝒜}}(\mathrm{\Phi }_\pm Q_a))Q_a)`$
$`=`$ $`\pm m_2^{\stackrel{~}{𝒜}}(m_0^{\stackrel{~}{𝒜}}(1)\mathrm{\Phi }_\pm )m_2^{\stackrel{~}{𝒜}}(\mathrm{\Phi }_\pm m_0^{\stackrel{~}{𝒜}}(1))\pm m_2^{\stackrel{~}{𝒜}}(m_1^{\stackrel{~}{𝒜}}(Q_b)\mathrm{\Phi }_\pm )`$
$`m_2^{\stackrel{~}{𝒜}}(\mathrm{\Phi }_\pm m_1^{\stackrel{~}{𝒜}}(Q_a))\pm m_2^{\stackrel{~}{𝒜}}(m_2^{\stackrel{~}{𝒜}}(Q_b^2)\mathrm{\Phi }_\pm )m_2^{\stackrel{~}{𝒜}}(\mathrm{\Phi }_\pm m_2^{\stackrel{~}{𝒜}}(Q_a^2))`$
$`=`$ $`\pm m_2^{\stackrel{~}{𝒜}}((m_0^{\stackrel{~}{𝒜}}(1)+m_1^{\stackrel{~}{𝒜}}(Q_b)+m_2^{\stackrel{~}{𝒜}}(Q_b^2)\mathrm{\Phi }_\pm )`$
$`m_2^{\stackrel{~}{𝒜}}(\mathrm{\Phi }_\pm (m_0^{\stackrel{~}{𝒜}}(1+m_1^{\stackrel{~}{𝒜}}(Q_a)+m_2^{\stackrel{~}{𝒜}}(Q_a^2))`$
$`=`$ $`0.`$
###### Lemma 2.14.
For $`\mathrm{\Phi }Tw(𝒜)(\alpha _0,\alpha _1)`$ and $`\mathrm{\Psi }Tw(𝒜)(\alpha _1,\alpha _2)`$, we have
$$m_1^{Tw(𝒜)}(m_2^{Tw(𝒜)}(\mathrm{\Psi }\mathrm{\Phi }))=(1)^{|\mathrm{\Phi }|}m_2^{Tw(𝒜)}(m_1^{Tw(𝒜)}(\mathrm{\Psi })\mathrm{\Phi })m_2^{Tw(𝒜)}(\mathrm{\Psi }m_1^{Tw(𝒜)}(\mathrm{\Phi })).$$
(2.26)
###### Proof.
$`m_1^{Tw(𝒜)}(m_2^{Tw(𝒜)}(\mathrm{\Psi }\mathrm{\Phi }))`$
$`=`$ $`m_1^{\stackrel{~}{𝒜}}(m_2(\mathrm{\Psi }\mathrm{\Phi }))+m_2^{\stackrel{~}{𝒜}}(Q_{a_2}m_2^{\stackrel{~}{𝒜}}(\mathrm{\Psi }\mathrm{\Phi }))+m_2^{\stackrel{~}{𝒜}}(m_2^{\stackrel{~}{𝒜}}(\mathrm{\Psi }\mathrm{\Phi })Q_{a_0})`$
$`=`$ $`(1)^{|\mathrm{\Phi }|}m_2^{\stackrel{~}{𝒜}}(m_1^{\stackrel{~}{𝒜}}(\mathrm{\Psi })\mathrm{\Phi }))m_2^{\stackrel{~}{𝒜}}(\mathrm{\Psi }m_1^{\stackrel{~}{𝒜}}(\mathrm{\Phi }))`$
$`(1)^{|\mathrm{\Phi }|}m_2^{\stackrel{~}{𝒜}}(m_2^{\stackrel{~}{𝒜}}(Q_{a_2}\mathrm{\Psi })\mathrm{\Phi })m_2^{\stackrel{~}{𝒜}}(\mathrm{\Psi }m_2^{\stackrel{~}{𝒜}}(\mathrm{\Phi }Q_{a_0}))`$
$`=`$ $`(1)^{|\mathrm{\Phi }|}m_2^{Tw(𝒜)}(m_1^{Tw(𝒜)}(\mathrm{\Psi })\mathrm{\Phi })(1)^{|\mathrm{\Phi }|}m_2^{\stackrel{~}{𝒜}}(m_2^{\stackrel{~}{𝒜}}(\mathrm{\Psi }Q_{a_1})\mathrm{\Phi })`$
$`m_2^{Tw(𝒜)}(\mathrm{\Psi }m_1^{Tw(𝒜)}(\mathrm{\Phi }))+m_2^{\stackrel{~}{𝒜}}(\mathrm{\Psi }m_2^{\stackrel{~}{𝒜}}(Q_{a_1}\mathrm{\Phi }))`$
$`=`$ $`(1)^{|\mathrm{\Phi }|}m_2^{Tw(𝒜)}(m_1^{Tw(𝒜)}(\mathrm{\Psi })\mathrm{\Phi })m_2^{Tw(𝒜)}(\mathrm{\Psi }m_1^{Tw(𝒜)}(\mathrm{\Phi })).`$
By the above two Lemmas, we have the following.
###### Proposition 2.15.
Let $`𝒜`$ be a unital $``$-graded $`A_{\mathrm{}}`$-category of index $`h`$ with $`m_n^𝒜=0`$, $`n3`$. A collection $`Ob(Tw(𝒜))`$, $`Tw(𝒜)(\alpha ,\beta )`$ and $`(m_0^{Tw(𝒜)},m_1^{Tw(𝒜)},m_2^{Tw(𝒜)})`$ given by Definition 2.9 and Definition 2.12 determines a structure of a differential $`(_2`$-$`)`$graded category. We denote it by $`Tw(𝒜)`$. ∎
###### Remark 2.16.
Note that the condition that $`𝒜`$ is $``$-graded is not necessary for the above definition of the category $`Tw(𝒜)`$ and the category $`D^b(𝒜)`$ below. We need only the $`_2`$-grading.
###### Definition 2.17.
Let $`𝒜`$ be a unital $``$-graded $`A_{\mathrm{}}`$-category of index $`h`$ with $`m_n^𝒜=0`$, $`n3`$. We construct the category $`D^b(𝒜)`$ called the bounded derived category of $`𝒜`$ as follows. The set of objects is given by
$$Ob(D^b(𝒜)):=Ob(Tw(𝒜)),$$
(2.27)
and the set of homomorphisms is given by
$`\mathrm{Hom}_{D^b(𝒜)}(\alpha ,\beta ):=`$ $`\mathrm{Ker}(m_1^{Tw(𝒜)}:Tw(𝒜)(\alpha ,\beta )_+Tw(𝒜)(\alpha ,\beta )_{})`$ (2.28)
$`/\mathrm{Im}(m_1^{Tw(𝒜)}:Tw(𝒜)(\alpha ,\beta )_{}Tw(𝒜)(\alpha ,\beta )_+).`$ (2.29)
Let $`TTw(𝒜)(\alpha ,\beta )_+`$ be a $`m_1^{Tw(𝒜)}`$-closed homomorphism. We define a mapping cone $`C(T)`$ as an object
$$C(T):=(a[1]b;Q_{C(T)}),Q_{C(T)}:=\left(\begin{array}{cc}Q_{a[1]}& 0\\ T& Q_b\end{array}\right).$$
(2.30)
$`C(T)`$ is well-defined since the Maurer-Cartan equation (2.16) for $`Q_{C(T)}`$ is equivalent to the equation $`m_1^{Tw(𝒜)}(T)=0`$ and the Maurer-Cartan equation (2.16) for $`Q_a`$ and $`Q_b`$. Note also that there are natural closed morphisms
$$\beta C(T),C(T)\alpha .$$
We define an exact triangle in the category $`D^b(𝒜)`$ as a triangle of the form
$$\alpha \stackrel{T}{}\beta C(T)\alpha [1],$$
(2.31)
for some $`T\mathrm{Ker}(m_1^{Tw(𝒜)}:Tw(𝒜)(\alpha ,\beta )_+Tw(𝒜)(\alpha ,\beta )_{})`$.
###### Theorem 2.18.
The category $`D^b(𝒜)`$ endowed with a shift functor $`[1]`$ and the class of exact triangles defined above becomes a triangulated category.
###### Proof.
The proof is essentiality the same as the known results in the usual situation. See for example, \[BK\],\[GM\], \[KS\], \[O1\] and \[Se1\]. ∎
###### Remark 2.19.
The twice of the shift functor $`[2]:=[1]^2`$ is isomorphic to the identity functor in $`D^b(𝒜)`$.
We shall add more objects to $`D^b(𝒜)`$ following \[Se1\].
###### Definition 2.20.
Consider the category $`D^\pi (𝒜)`$ whose objects are pairs $`(X,p)`$ where $`XOb(D^b(𝒜))`$ and $`p\mathrm{Hom}_{D^b(𝒜)}(X,X)`$ an idempotent endomorphism, and whose spaces of homomorphisms are $`\mathrm{Hom}_{D^\pi (𝒜)}((X_0,p_0),(X_1,p_1)):=p_1\mathrm{Hom}_{D^b(𝒜)}(X_0,X_1)p_0`$. The category $`D^\pi (𝒜)`$ is called the split-closed derived category of $`𝒜`$.
It is known that $`D^\pi (𝒜)`$ is again a triangulated category (see \[BS\]).
###### Remark 2.21.
Since any projective module over the polynomial ring $`[x_1,\mathrm{},x_n]`$ is a free module, we have $`D^b(𝒜_f)D^\pi (𝒜_f)`$. Therefore, we shall study in this paper only the category $`D^b(𝒜_f)`$.
It is not difficult to see that our category $`D^b(𝒜_f)`$ is equivalent to the category of matrix factorizations, or equivalently, the category of maximal Cohen-Macaulay modules over $`[x_1,\mathrm{},x_n]/(f)`$ without free summands studied by Eisenbud \[E\], Knörrer \[K\], Buchweitz-Greuel-Schreyer \[BGS\], Orlov\[O1\] and other people $`(`$see the book by Yoshino \[Y\] for the details on maximal Cohen-Macaulay modules$`)`$. Indeed, we can construct a functor from the category of matrix factorizations to $`D^b(𝒜_f)`$ once we choose a basis of the free module over the polynomial ring. Note that any object isomorphic to a direct sum of the following objects
$$\left(\begin{array}{cc}0& f\\ 1& 0\end{array}\right),\left(\begin{array}{cc}0& 1\\ f& 0\end{array}\right)$$
becomes the zero object of the category $`D^b(𝒜_f)`$.
Next, we shall define the $``$-equivariant bounded derived category $`D_{}^b(𝒜)`$ of $`𝒜`$. Let $`\alpha :=(a;Q_a)`$ and $`\beta :=(b;Q_b)`$ be graded twisted complexes. We put
$$Tw_{}(𝒜)(\alpha ,\beta ):=\underset{q}{}Tw_{}^q(𝒜)(\alpha ,\beta ),$$
where
$$Tw_{}^q(𝒜)(\alpha ,\beta ):=\{\begin{array}{c}\stackrel{~}{𝒜}^q(a,b)_+,q2,\hfill \\ \stackrel{~}{𝒜}^q(a,b)_{},q12.\hfill \end{array}$$
(2.32)
Since $`Tw_{}^q(𝒜)(\alpha ,\beta )Tw(𝒜)(\alpha ,\beta )`$, we can define a $`k`$-multilinear maps $`m_n^{Tw_{}(𝒜)}`$ by restricting $`m_n^{Tw(𝒜)}`$ to the subspaces.
###### Proposition 2.22.
Let $`𝒜`$ be a unital $``$-graded $`A_{\mathrm{}}`$-category of index $`h`$ with $`m_n^𝒜=0`$, $`n3`$. A collection $`Ob(Tw_{}(𝒜))`$, $`Tw_{}(\alpha ,\beta )`$ and $`m_n^{Tw_{}(𝒜)}`$ given above determines a $``$-graded $`A_{\mathrm{}}`$-category with $`m_n0`$ only if $`n=1,2`$, i.e., a differential graded $`(`$DG$`)`$ category in the usual sense. We denote it by $`Tw_{}(𝒜)`$. ∎
###### Definition 2.23.
Let $`𝒜`$ be a unital $``$-graded $`A_{\mathrm{}}`$-category of index $`h`$ with $`m_n^𝒜=0`$, $`n3`$. We call the cohomology category of $`Tw_{}(𝒜)`$ the $``$-equivariant bounded derived category of $`𝒜`$ and denote by $`D_{}^b(𝒜)`$. More precisely, the set of objects is given by
$$Ob(D_{}^b(𝒜)):=Ob(Tw_{}(𝒜)),$$
(2.33)
and the set of homomorphisms is given by
$`\mathrm{Hom}_{D_{}^b(𝒜)}(\alpha ,\beta ):=`$ $`\mathrm{Ker}(m_1^{Tw_{}(𝒜)}:Tw_{}(𝒜)^0(\alpha ,\beta )Tw_{}(𝒜)^1(\alpha ,\beta ))`$ (2.34)
$`/\mathrm{Im}(m_1^{Tw_{}(𝒜)}:Tw_{}(𝒜)^1(\alpha ,\beta )Tw_{}(𝒜)^0(\alpha ,\beta )).`$ (2.35)
Let $`TTw_{}(𝒜)^0(\alpha ,\beta )`$ be a $`m_1^{Tw_{}(𝒜)}`$-closed homomorphism. As in the case for $`D^b(𝒜)`$, we define a mapping cone $`C(T)`$ as an object
$$C(T):=(a[1]b;Q_{C(T)}),Q_{C(T)}:=\left(\begin{array}{cc}Q_{a[1]}& 0\\ T& Q_b\end{array}\right).$$
(2.36)
We define an exact triangle in the category $`D_{}^b(𝒜)`$ as a triangle of the form
$$\alpha \stackrel{T}{}\beta C(T)\alpha [1],$$
(2.37)
for some $`T\mathrm{Ker}(m_1^{Tw_{}(𝒜)}:Tw_{}(𝒜)^0(\alpha ,\beta )Tw_{}(𝒜)^1(\alpha ,\beta ))`$.
###### Theorem 2.24.
The category $`D_{}^b(𝒜)`$ endowed with a shift functor $`[1]`$ and the class of exact triangles defined above becomes a triangulated category.
###### Proof.
As in the case for $`D^b(𝒜)`$, the proof is essentiality the same as the known results in the usual situation. ∎
###### Remark 2.25.
The twice of the shift functor $`[2]`$ is not isomorphic to the identity functor in $`D_{}^b(𝒜)`$.
Consider the differential graded (DG) functor $`\mathrm{Tot}`$ as in \[BK\]
$$\mathrm{Tot}:Tw_{}(Tw_{}(𝒜))Tw_{}(𝒜),(\underset{i=1}{\overset{k}{}}(a_i;Q_{a_i});T)(\underset{i=1}{\overset{k}{}}a_i;Q+T),$$
(2.38)
where
$$Q:=\left(\begin{array}{ccc}Q_{a_1}& 0& 0\\ 0& \mathrm{}& 0\\ 0& 0& Q_{a_k}& \end{array}\right),$$
(2.39)
and $`TTw_{}(𝒜)(_{i=1}^k(a_i;Q_{a_i}),_{i=1}^k(a_i;Q_{a_i}))`$ $`(\stackrel{~}{𝒜}(_{i=1}^ka_i,_{i=1}^ka_i)`$) satisfies
$$m_1^{Tw_{}(𝒜)}(T)+m_2^{Tw_{}(𝒜)}(T^2)=0.$$
It is well-defined since
$`m_0^{\stackrel{~}{𝒜}}(1)+m_1^{\stackrel{~}{𝒜}}(Q+T)+m_2^{\stackrel{~}{𝒜}}((Q+T)^2)`$
$`=`$ $`m_0^{\stackrel{~}{𝒜}}(1)+{\displaystyle \underset{i=1}{\overset{k}{}}}m_1^{\stackrel{~}{𝒜}}(Q_{a_i})+{\displaystyle \underset{i=1}{\overset{k}{}}}m_2^{\stackrel{~}{𝒜}}(Q_{a_i}^2)+m_1^{\stackrel{~}{𝒜}}(T)+m_2^{\stackrel{~}{𝒜}}(TQ)+m_2^{\stackrel{~}{𝒜}}(QT)+m_2^{\stackrel{~}{𝒜}}(T^2)`$
$`=`$ $`m_0^{\stackrel{~}{𝒜}}(1)+{\displaystyle \underset{i=1}{\overset{k}{}}}m_1^{\stackrel{~}{𝒜}}(Q_{a_i})+{\displaystyle \underset{i=1}{\overset{k}{}}}m_2^{\stackrel{~}{𝒜}}(Q_{a_i}^2)+m_1^{Tw_{}(𝒜)}(T)+m_2^{Tw_{}(𝒜)}(T^2).`$
Now the following statement is easily shown as in \[BK\] where they consider the case when $`𝒜`$ is a DG category, i.e., the case when $`m_0^𝒜=0`$ and $`h=1`$ in our terminology.
###### Proposition 2.26.
$`\mathrm{Tot}`$ is an equivalence of DG categories. ∎
Indeed, by definition of the twisted complexes and the differential on them, we see that $`(Tw_{}(Tw_{}(𝒜))(A,B),m_1^{Tw_{}(Tw_{}(𝒜))})`$ and $`(Tw_{}(𝒜)(\mathrm{Tot}(A),\mathrm{Tot}(B)),m_1^{Tw_{}(𝒜)})`$ are the same as complexes.
###### Corollary 2.27.
$`D_{}^b(𝒜)`$ is an enhanced triangulated category in the sense of Bondal-Kapranov \[BK\]. ∎
Let us consider our category $`D_{}^b(𝒜_f)`$ a little bit in detail.
###### Example 2.28.
Let $`\alpha :=(a=a_+a_{}[1];Q_a)`$ and $`\beta :=(b=b_+b_{}[1]:Q_b)`$ be objects of $`Tw_{}(𝒜_f)`$. Then the space of homomorphisms is of the following form:
$$\mathrm{\Phi }Tw_{}^q(\alpha ,\beta ),q2\mathrm{\Phi }=\left(\begin{array}{cc}\mathrm{\Phi }_{++}& 0\\ 0& \mathrm{\Phi }_{}\end{array}\right),\mathrm{\Phi }_{\pm \pm }\stackrel{~}{𝒜_f}^q(a_\pm ,b_\pm )_+,$$
$$\mathrm{\Phi }Tw_{}^q(\alpha ,\beta ),q12\mathrm{\Phi }=\left(\begin{array}{cc}0& \mathrm{\Phi }_+\\ \mathrm{\Phi }_+& 0\end{array}\right),\mathrm{\Phi }_\pm \stackrel{~}{𝒜_f}^{q1}(a_\pm ,b_{})_+,$$
and the coboundary operator $`m_1^{Tw(𝒜_f)}`$ becomes the differential of the usual form
$$Q_b\mathrm{\Phi }(1)^q\mathrm{\Phi }Q_a,\mathrm{\Phi }Tw_{}^q(\alpha ,\beta ).$$
Note that if $`\mathrm{\Phi }\stackrel{~}{𝒜_f}^q(\alpha ,\beta )`$, then
$$E\mathrm{\Phi }R_\beta \mathrm{\Phi }+\mathrm{\Phi }R_\alpha =q\mathrm{\Phi },$$
(2.40)
where we put
$$R_\alpha :=\mathrm{diag}(\frac{2k_1}{h},\mathrm{},\frac{2k_m}{h},\frac{2l_1}{h}1,\mathrm{},\frac{2l_m}{h}1),a=\underset{i=1}{\overset{m}{}}a\{\frac{2k_i}{h}\}\underset{i=1}{\overset{m}{}}a\{\frac{2l_i}{h}\}[1],$$
and
$$R_\beta :=\mathrm{diag}(\frac{2k_1^{}}{h},\mathrm{},\frac{2k_m^{}^{}}{h},\frac{2l_1^{}}{h}1,\mathrm{},\frac{2l_m^{}^{}}{h}1),b=\underset{i=1}{\overset{m^{}}{}}a\{\frac{2k_i^{}}{h}\}\underset{i=1}{\overset{m^{}}{}}a\{\frac{2l_i^{}}{h}\}[1].$$
By integrating the equation (2.40), we get for $`\lambda `$,
$$e^{\lambda R_\beta }\mathrm{\Phi }(e^{\lambda \frac{2a_1}{h}}x_1,\mathrm{},e^{\lambda \frac{2a_n}{h}}x_n)e^{\lambda R_\alpha }=e^{q\lambda }\mathrm{\Phi }(x_1,\mathrm{},x_n).$$
(2.41)
This is the analogue of the homogeneity condition discussed in \[HW\].
Consider the $``$-action defined by $`x_i\mathrm{exp}(2\pi \sqrt{1}pa_i/h)x_i`$, $`p`$. It is clear that $`f`$ is invariant under this $``$-action. Note also that $`\{\frac{2k}{h}\}`$, $`k`$ can be considered as the irreducible representations of $``$. For a graded twisted complex $`\alpha :=(a;Q_a)`$, put
$$S_\alpha :=\mathrm{diag}(\frac{2k_1}{h},\mathrm{},\frac{2k_m}{h},\frac{2l_1}{h},\mathrm{},\frac{2l_m}{h}),a=\underset{i=1}{\overset{m}{}}a\{\frac{2k_i}{h}\}\underset{i=1}{\overset{m}{}}a\{\frac{2l_i}{h}\}[1].$$
Since $`Q_a\stackrel{~}{𝒜}^1(a,a)_{}`$, the similar equation as (2.41) shows that $`Q_a`$ is equivariant with respect to the $``$-action, i.e., we have
$$e^{\pi \sqrt{1}S_\alpha }Q_a(e^{\frac{2\pi \sqrt{1}a_1}{h}}x_1,\mathrm{},e^{\frac{2\pi \sqrt{1}a_n}{h}}x_n)e^{\pi \sqrt{1}S_\alpha }=Q_a(x_1,\mathrm{},x_n).$$
(2.42)
One can show that there is also the $``$-action on the space of homomorphisms by (2.41). For $`\mathrm{\Phi }_\pm \stackrel{~}{𝒜_f}^q(\alpha ,\beta )_\pm `$, we have
$$e^{\pi \sqrt{1}S_\beta }\mathrm{\Phi }_\pm (e^{\frac{2\pi \sqrt{1}a_1}{h}}x_1,\mathrm{},e^{\frac{2\pi \sqrt{1}a_n}{h}}x_n)e^{\pi \sqrt{1}S_\alpha }=\pm e^{\varphi \sqrt{1}q}\mathrm{\Phi }_\pm (x_1,\mathrm{},x_n).$$
(2.43)
Therefore, if $`\mathrm{\Phi }`$ is even $`(`$odd$`)`$, then $`\mathrm{\Phi }`$ is $``$-invariant if and only if $`q2`$ $`(q12)`$. These facts lead us to our definition of $``$-equivariant derived category $`D_{}^b(𝒜_f)`$ of $`𝒜_f`$.
Note that the above $``$-action on $`𝒜_f`$ factors through $`/h`$. The category whose set of objects is the set of $`/h`$-equivariant matrix factorizations and the space of morphisms is $`/h`$-invariant homomorphisms between matrix factorizations are called in physics the category of D-branes in Landau-Ginzburg $`(/h`$-$`)`$orbifolds $`(`$see for example \[HW\]$`)`$ We can construct it by considering the $`/h`$-equivariant version of $`D^b(𝒜_f)`$. Indeed, we can show that it is equivalent to $`D_{}^b(𝒜_f)/[2]`$. In order to recover the $``$-grading by the shift functor, we introduced here the translation $`\{2/h\}`$ and defined a new category $`D_{}^b(𝒜_f)`$.
## 3. $`D_{}^b(𝒜_f)`$ and representations of Dynkin quivers
The following is our main theorem in this paper.
###### Theorem 3.1.
Let us put $`f(x):=x^h[x]`$ for $`h2`$ and consider the unital $``$-graded $`A_{\mathrm{}}`$-category $`𝒜_f`$ of index $`h`$. Then we have the following equivalence of triangulated categories
$$D_{}^b(𝒜_f)D^b(\mathrm{mod}B),$$
(3.1)
where $`B`$ is the path algebra of the following Dynkin quiver of type $`A_{h1}:`$
$$_1_2\mathrm{}_{h2}_{h1},$$
(3.2)
$`(`$the algebra of upper triangular matrices over $`k)`$, and $`D^b(\mathrm{mod}B)`$ is the bounded derived category of finitely generated right $`B`$-modules.
###### Remark 3.2.
The above equivalence for $`h=2`$, $`D_{}^b(𝒜_{x^2})D^b(\mathrm{mod})`$, gives the simplest example of Knörrer’s periodicity.
###### Proof.
It is not difficult to see that our category $`D_{}^b(𝒜_f)`$ is a Krull-Schmidt category, the spaces of homomorphisms are finite dimensional and the endomorphism rings of indecomposable objects are local rings. See, for example, section 5 of \[KR\] for the proof of the general $`f`$ which defines an isolated singularity. Therefore, we first study the set of isomorphism classes of indecomposable objects. We use the fact that the Auslander-Reiten quiver of the category $`D^b(𝒜_f)`$ of matrix factorizations for $`f`$ is given by
$$[Q_1][Q_2]\mathrm{}[Q_{h2}][Q_{h1}],$$
(3.3)
where
$$Q_l=\left(\begin{array}{cc}0& x^{hl}\\ x^l& 0\end{array}\right),l=1,\mathrm{},h1,i,$$
and each morphism corresponding to the arrow from left to right is given by $`\mathrm{diag}(1,x)`$ and the one from right to left is given by $`\mathrm{diag}(x,1)`$. See \[AR\] and also \[O1\]. Hence we have the following.
###### Lemma 3.3.
The set of isomorphism classes of all indecomposable objects of $`D_{}^b(𝒜_f)`$ is given by
$$\left\{[M_{l,i}],l=1,\mathrm{},h1,i\right\},$$
(3.4)
where
$$M_{l,i}:=(a\{\frac{2i}{h}\}a\{\frac{2(l+i)}{h}\}[1];Q_l).$$
(3.5)
We also have
$$\left(\begin{array}{cc}1& 0\\ 0& x\end{array}\right)\mathrm{Hom}_{D_{}^b(𝒜_f)}(M_{l,i},M_{l+1,i}),\left(\begin{array}{cc}x& 0\\ 0& 1\end{array}\right)\mathrm{Hom}_{D_{}^b(𝒜_f)}(M_{l,i},M_{l1,i+1}),$$
(3.6)
and hence
$$\mathrm{Hom}_{D_{}^b(𝒜_f)}(M_{k,i},M_{l,j})0,\text{only if }k+2il+2j.$$
(3.7)
In particular,
$$\mathrm{Hom}_{D_{}^b(𝒜_f)}(M_{k,0},M_{l,0})=\{\begin{array}{c},\mathrm{if}kl,\hfill \\ 0,\mathrm{if}k>l.\hfill \end{array}$$
(3.8)
###### Proof.
One can easily show by direct computations. ∎
###### Remark 3.4.
Note that $`M_{l,i}[1]M_{hl,l+i}`$.
Serre duality holds in our category $`D_{}^b(𝒜_f)`$.
###### Lemma 3.5.
There are isomorphisms as $``$-vector spaces
$$\mathrm{Hom}_{D_{}^b(𝒜_f)}(M_{k,i},M_{l,j})\mathrm{Hom}_{D_{}^b(𝒜_f)}(M_{l,j},M_{k,i1}[1])^{},\text{for all}1k,lh1,i,j.$$
(3.9)
###### Proof.
There is a trace map \[KL2\]
$$Tr_k:\stackrel{~}{𝒜_f}^{1\frac{2}{h}}(M_{k,i},M_{k,i})_{},\mathrm{\Phi }\frac{1}{h1}\mathrm{Res}\left[\frac{Str(dQ_k\mathrm{\Phi })}{\frac{f}{x}}\right],$$
(3.10)
where
$$Str(dQ_k\mathrm{\Phi }):=\left[(hk)x^{hk1}\mathrm{\Phi }_+kx^{k1}\mathrm{\Phi }_+\right]dx,\mathrm{\Phi }=\left(\begin{array}{cc}0& \mathrm{\Phi }_+\\ \mathrm{\Phi }_+& 0\end{array}\right).$$
$`Tr_k(\mathrm{\Phi })=0`$ if $`\mathrm{\Phi }=Q_k\mathrm{\Psi }+\mathrm{\Psi }Q_k`$ for some $`\mathrm{\Psi }`$ since $`dQ_kQ_k+Q_kdQ_k=df1_{2\times 2}`$.
Note that
$$\mathrm{\Phi }_k:=\left(\begin{array}{cc}0& x^{hk1}\\ x^{k1}& 0\end{array}\right)\stackrel{~}{𝒜_f}^{1\frac{2}{h}}(M_{k,i},M_{k,i})_{},Tr_k(\mathrm{\Phi }_k)=1,Q_k\mathrm{\Phi }_k+\mathrm{\Phi }_kQ_k=0.$$
Therefore, under the isomorphism given by
$$\stackrel{~}{𝒜_f}^{1\frac{2}{h}}(M_{k,i},M_{k,i})_{}\stackrel{~}{𝒜_f}^0(M_{k,i},M_{k,i1}[1])_+,\left(\begin{array}{cc}0& \mathrm{\Phi }_+\\ \mathrm{\Phi }_+& 0\end{array}\right)\left(\begin{array}{cc}\mathrm{\Phi }_+& 0\\ 0& \mathrm{\Phi }_+\end{array}\right),$$
$`\mathrm{\Phi }_k`$ determines an element of $`\mathrm{Hom}_{D_{}^b(𝒜_f)}(M_{l,j},M_{k,i1}[1])`$. Moreover, from the knowledge of the Auslander-Reiten quiver (3.3) of $`D^b(𝒜_f)`$, we see that $`\mathrm{Hom}_{D_{}^b(𝒜_f)}(M_{l,j},M_{k,i1}[1])\mathrm{\Phi }_k`$.
It is not difficult to see that the following pairings
$$\stackrel{~}{𝒜_f}^{\frac{2m}{h}}(M_{k,i},M_{l,j})_+\stackrel{~}{𝒜_f}^{1\frac{2}{h}\frac{2m}{h}}(M_{l,j},M_{k,i})_{}\stackrel{~}{𝒜_f}^{1\frac{2}{h}\frac{2m}{h}}(M_{k,i},M_{k,i})_{}\stackrel{Tr_k}{},m,$$
induce the perfect pairings
$$\mathrm{Hom}_{D_{}^b(𝒜_f)}(M_{k,i},M_{l,j})\mathrm{Hom}_{D_{}^b(𝒜_f)}(M_{l,j},M_{k,i1}[1])\mathrm{Hom}_{D_{}^b(𝒜_f)}(M_{k,i},M_{k,i1}[1])\stackrel{Tr_k}{}.$$
###### Remark 3.6.
$`S:=\{\frac{2}{h}\}[1]`$ is the Serre functor on $`D_{}^b(𝒜_f)`$. In particular, we have $`S^h=[h2]`$. Therefore $`D_{}^b(𝒜_f)`$ is a fractional noncommutative Calabi-Yau manifold of dimension $`12/h`$ in the sense of \[So\].
Combining the above Serre duality and the data of the Auslander-Reiten quiver (3.3) of $`D^b(𝒜_f)`$, we see that there are no higher extensions among $`\{M_{l,0}\}`$.
###### Corollary 3.7.
For $`m0`$, we have
$$\mathrm{Hom}_{D_{}^b(𝒜)}(M_{k,0},M_{l,0}[m])=0,\text{for all}i,j=1,\mathrm{},h1.$$
###### Corollary 3.8.
$`D^b(\mathrm{mod}B)`$ is a full triangulated subcategory of $`D_{}^b(𝒜_f)`$.
###### Proof.
Use the fact that $`(M_{1,0},\mathrm{},M_{h1,0})`$ is a strongly exceptional collection and
$$B\underset{i,j=1}{\overset{h1}{}}\mathrm{Hom}_{D_{}^b(𝒜)}(M_{k,0},M_{l,0}).$$
Since $`D_{}^b(𝒜_f)`$ is an enhanced triangulated category, we can apply the theorem by Bondal-Kapranov $`(`$\[BK\] Theorem 1$`)`$. ∎
Note that the number of indecomposable objects of $`D_{}^b(𝒜_f)/[2]`$ is
$$\mathrm{\#}\left\{[M_{l,i}]\right|l=1,\mathrm{},h1,i/h\}=(h1)h,$$
which is the number of roots for the root system $`A_{h1}`$. This number coincides with the number of indecomposable objects of $`D^b(\mathrm{mod}B)/[2]`$ by Gabriel’s theorem \[G\]. Therefore, $`D^b(\mathrm{mod}B)/[2]D_{}^b(𝒜_f)/[2]`$. This proves Theorem 3.1. ∎
###### Remark 3.9.
The similar proof can be applied for $`D_n`$ and $`E_6,E_7,E_8`$ cases since the heart of our proof is to use the Auslander-Reiten quivers of $`D^b(𝒜_f)`$, the fact that any matrix factorization over ADE singularities is gradable, the Serre duality and the theorem by Gabriel on the number of indecomposables. They are well-known or can be shown by direct calculations with explicit presentations of matrix factorizations. We shall discuss this in detail in the next paper \[KST\].
## 4. Stability condition on $`D_{}^b(𝒜_f)`$
In this section, we will briefly discuss on a stability condition on $`D_{}^b(𝒜_f)`$.
###### Definition 4.1.
Let $`\alpha :=(_{i=1}^na\{\frac{2k_i}{h}\}_{i=1}^na\{\frac{2l_i}{h}\}[1];Q_a)`$ be an object of $`D_{}^b(𝒜_f)`$ such that $`Q_a`$ is reduced, i.e., each matrix element of $`Q_a`$ is in the maximal ideal generated by $`(x_1,\mathrm{},x_n)`$. Then we call the real number
$$\varphi _\alpha :=\frac{1}{2n}\mathrm{Tr}S_a\frac{1}{2},S_a:=\mathrm{diag}(\frac{2k_1}{h},\mathrm{},\frac{2k_n}{h},\frac{2l_1}{h},\mathrm{},\frac{2l_n}{h})$$
(4.1)
phase of the object $`\alpha `$.
###### Example 4.2.
Let $`f:=x^{n+1}`$ and consider the objects
$$M_{l,i}:=(a\{\frac{2i}{h}\}a\{\frac{2(l+i)}{h}\}[1];\left(\begin{array}{cc}0& x^{hl}\\ x^l& 0\end{array}\right)).$$
(4.2)
Then
$$\varphi _{M_{l,i}}=\frac{l+2i}{h}\frac{1}{2}.$$
###### Definition 4.3.
Let $`\omega :=\mathrm{exp}2\pi \sqrt{1}/h`$. For $`\alpha =(_{i=1}^na\{\frac{2k_i}{h}\}_{i=1}^na\{\frac{2l_i}{h}\}[1];Q_a)`$, we define a $``$-linear map $`Z_\omega :K_0(D_{}^b(𝒜_f))`$ as follows$`:`$
$$Z_\omega (\alpha ):=\underset{i=1}{\overset{n}{}}(\omega ^{k_i}\omega ^{l_i}).$$
(4.3)
By the above example, we see that $`\varphi _{M_{l,i}}`$ is the phase of $`Z_\omega ([M_{l,i}])`$:
###### Proposition 4.4.
$$Z_\omega ([M_{l,i}])=2\mathrm{sin}(\frac{l}{h}\pi )e^{\pi \sqrt{1}\varphi _{M_{l,i}}}.$$
(4.4)
Since we know that all indecomposable objects in $`D_{}^b(𝒜_f)`$ for $`f=x^{n+1}`$ have definite phases, we can define a stability condition on $`D_{}^b(𝒜_f)`$.
###### Theorem 4.5.
Let $`f:=x^{n+1}`$ and $`P(\varphi )`$ be the full additive subcategory of $`D_{}^b(𝒜_f)`$ whose objects have phase $`\varphi `$. Then $`P(\varphi )`$ and $`Z_\omega `$ define a stability condition on $`D_{}^b(𝒜_f)`$ in the sense of Bridgeland \[B\].
More precisely, $`P(\varphi )`$ and $`Z_\omega `$ satisfy the following properties$`:`$
1. if $`MP(\varphi )`$, then $`Z_\omega (M)=m(M)\mathrm{exp}(\sqrt{1}\pi \varphi )`$ for some $`m(M)_0`$,
2. for all $`\varphi ,`$ $`P(\varphi +1)=P(\varphi )[1]`$,
3. if $`\varphi _1>\varphi _2`$ and $`M_iP(\varphi _i)`$, then $`\mathrm{Hom}_{D_{}^b(𝒜_f)}(M_1,M_2)=0`$,
4. for each nonzero object $`MD_{}^b(𝒜_f)`$, there is a finite sequence of real numbers
$$\varphi _1>\varphi _2>\mathrm{}>\varphi _n$$
and a collection of exact triangles
$$M_{i1}M_iN_iM_{i1}[1],M_n:=M,M_0:=0$$
with $`N_jP(\varphi _j)`$ for all $`j`$.
The space of stability conditions for $`D_{}^b(𝒜_f)`$ should be isomorphic to the base space of the universal unfolding of $`f`$ by the mirror symmetry. Therefore we expect that there exists a natural Frobenius (K. Saito’s flat) structure on the space of stability conditions and the stability condition constructed above should correspond to the origin of the base space of the universal unfolding. We shall study this in detail elsewhere.
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# Nonlinearity induced destruction of resonant tunneling in the Wannier-Stark problem
## Abstract
We present detailed numerical results on the dynamics of a Bose-Einstein condensate in a tilted periodic optical lattice over many Bloch periods. We show that an increasing atom-atom interaction systematically affects coherent tunneling, and eventually destroys the resonant tunneling peaks.
Experiments with cold and ultracold atoms made it possible in the last decade to prepare and control the centre-of-mass motion of atoms with unprecedented precision. Many toy models of either many-body solid state physics BOfirst ; BECBOfirst ; pisaBO ; nist ; ingu or of simple Hamiltonian systems, whose complexity arises from an external driving force driven , were realized with the exceptional control offered by static or time-dependent optical potentials.
Particularly Bose-Einstein condensates (BEC) whose initial momentum spread can be adjusted in width and absolute position have proved to be an extremely helpful experimental tool pisaBO ; nist ; fallani ; pisaASY ; instabil . In addition, a BEC offers interesting new features originating from the intrinsic interactions between the atoms. Examples of such effects are new quantum phases greiner , soliton-like motion eiermann , the occurrence of energetic or dynamical instabilities in condensates fallani ; instabil ; wuniu ; zheng , or the decay and subsequent revival of Bloch oscillations (BO) revival .
We focus on the evolution of a BEC loaded into a one-dimensional lattice and subjected to an additional static force $`F`$, which is most easily realized and controlled by accelerating the optical lattice BOfirst ; pisaBO . In previous experiments, a BEC was accelerated to allow for a single crossing of the Brillouin zone (BZ), and two effects were observed: for large accelerations, an enhanced tunneling probability from the ground state band to the first excited band due to the atom-atom interaction was measured pisaBO ; pisaASY . Secondly, for smaller accelerations (where tunneling is negligible) signatures of a dynamical instability in the BEC were observed instabil ; footnote . By contrast, here we investigate the dynamics of a BEC performing many Bloch oscillations (BO), and we ask ourselves how the atom-atom interaction affects tunneling for a sequence of BZ crossings. In particular, we scan $`F`$ to study the impact of the atom-atom interaction on resonantly enhanced tunneling (RET), for which the standard Landau-Zener prediction is modified even in the absence of interactions kolo . The RET leads to a faster decay of the Wannier states trapped in the potential wells. With the survival probability and the recurrence probability (see Eqs. (4) and (6) below) we present two consistent measures for the nonlinear RET which define clear experimental signatures of the destruction of the coherent tunneling process inside the periodic potential.
If we neglect interactions for a moment, our system will be described by the Hamiltonian
$$H=\frac{\mathrm{}^2}{2M}\frac{d^2}{dx^2}+V\mathrm{sin}^2\left(\frac{\pi x}{d_L}\right)+Fx.$$
(1)
Here $`d_L`$ is the spatial period of the optical lattice with maximal amplitude $`V`$, and $`M`$ the atomic mass. Eq. (1) defines the well-known Wannier-Stark problem, which gives rise to BO with period $`T_{\mathrm{Bloch}}=h/d_LF`$ ($`h`$ is Planck’s constant). If tunneling is small, we can view the system as moving at a constant speed in momentum space within the fundamental BZ. At the zone edge, most of the wave packet is reflected (giving rise to BO) while a small part can tunnel across the first band gap to the next higher-lying energy band and then escape quickly by successive tunneling events across the smaller (higher) band gaps. Landau-Zener theory predicts a decay rate kolo
$$\mathrm{\Gamma }(F)Fe^{\frac{b}{F}},$$
(2)
where $`b`$ is proportional to the square of the energy gaps. Eq. (2) is modified by RET which occurs when two Wannier-Stark levels in neighboring potential wells are coupled strongly due to their accidental degeneracy. The RET results in pronounced peaks in the tunneling rates, e.g., as a function of $`1/F`$, on top of the global exponential decay described by (2) kolo . In this paper we investigate the impact of the effective shift of the Wannier-Stark levels by a nonlinear interaction term.
For the linear problem (1), the decay rates have been measured previously in the regime of short life-times in the ground state band (of the order of 100 $`\mu `$s), where $`\mathrm{\Gamma }(F)`$ is essentially smooth Raizen1997 . Since RET is a coherent quantum effect, the peaks should be sensitively affected by the atom-atom interaction, which can be varied experimentally by changing either the density of the BEC or through the atom-atom scattering potential via a Feshbach resonance wieman2001 . Our results are a consequence of many sequential Landau-Zener events, and they show the destruction of a RET peak with increasing interaction strength, in a regime which is experimentally accessible.
We use a fully 3D Gross-Pitaevskii equation (GPE) Mannella1998 to describe the temporal evolution of a BEC which is subject to realistic potentials:
$`\text{i}\mathrm{}{\displaystyle \frac{}{t}}\psi (\stackrel{}{r},t)=[{\displaystyle \frac{\mathrm{}^2}{2M}}^2+{\displaystyle \frac{1}{2}}M(\omega _x^2x^2+\omega _\mathrm{r}^2\rho ^2)+`$
$`V\mathrm{sin}^2\left({\displaystyle \frac{\pi x}{d_L}}\right)+Fx+gN\left|\psi (\stackrel{}{r},t)|^2\right]\psi (\stackrel{}{r},t).`$ (3)
$`\psi (\stackrel{}{r},t)`$ represents the condensate wave function, and the frequencies $`\omega _x`$ and $`\omega _\mathrm{r}`$ characterize the longitudinal and transverse harmonic confinement (here with cylindrical symmetry: $`\rho =\sqrt{y^2+z^2}`$). We fixed $`d_L=1.56\mu \mathrm{m}`$ and $`V/E_\mathrm{R}=5`$ for our computations, with the recoil energy $`E_\mathrm{R}=p_\mathrm{R}^2/2M`$ for $`p_\mathrm{R}=\mathrm{}\pi /d_L`$, and the recoil period $`T_\mathrm{R}=h/E_\mathrm{R}`$. The above values for $`d_L`$ and $`V`$ were realized in the experiments reported in pisaBO ; pisaASY ; instabil based on two laser beams propagating at an angle different from $`\pi `$. In Eq. (3), the nonlinear coupling constant is given by $`g=4\pi \mathrm{}^2a_s/M`$, where $`a_s`$ is the $`s`$-wave scattering length and $`N`$ the number of atoms in the BEC O1998 ; Mannella1998 . The dimensionless nonlinearity $`C=gn_0/(8E_\mathrm{R})`$ is computed from the peak density of the initial state of the condensate, with $`C=0.027\mathrm{}0.31`$ for the experimentally investigated range of pisaBO , and with $`C=0.5`$ reached in chu . Here we focus on $`C>0`$, but report briefly also on attractive interactions with $`C<0`$. The latter case leads to a fundamentally different behavior of the system because the collapse of the condensate introduces an additional time scale, which for experimentally relevant parameters is of the order of 10 msec WHS2005 ; wieman2001 (slightly longer than $`T_{\mathrm{Bloch}}=1.8\mathrm{}3.0\mathrm{msec}`$ here).
The GPE (3) is numerically integrated using finite difference propagation, adapted by a predictor-corrector estimate to reliably evaluate the nonlinear interaction Mannella1998 . Since our system is essentially the problem of a constantly accelerated particle for the part of the wave function which has tunnelled out of the first BZ already, one must be careful with the application of absorbing boundary conditions or complex coordinate methods GPabs ; SP2004 . To avoid any spurious effects due to the fast spreading, we use a large numerical basis. In this way, we fully cover the 3D expansion of the entire wave packet, including its tunnelled tail, without the use of non-Hermitian potentials. The initial state propagated by (3) is the relaxed condensate wave function, adiabatically loaded into the confining potential given by the harmonic trap and the optical lattice (with $`F=0`$). Approximate analytic forms of the relaxed state are found, e.g., in pedri , but we used an imaginary time propagation to reliably compute the initial state for $`C>0`$.
The linear decay rates for non-interacting atoms in the optical lattice are computed from the spectrum of the 1D Wannier-Stark problem of Eq. (1) using, e.g., the method of kolo . Those linear rates are plotted in Fig. 1. The maxima in the rates occur when $`Fd_Lm`$ (with $`m`$ integer) is close to the difference between the first two energy bands (averaged over the BZ) of the $`F=0`$ problem kolo . The actual peaks are slightly shifted with respect to the above estimate (marked by arrows in the inset of Fig. 1), owing to a field-induced level shift kolo .
Experimentally, the most easily measurable quantity is the momentum distribution of the BEC obtained from a free expansion after the evolution inside the lattice. From the momentum distribution we determine the survival probability by projection of the evolved state $`\psi (\stackrel{}{p},t)`$ onto the support of the initial state
$$P_{\mathrm{sur}}(t)_{p_\mathrm{c}}^{p_\mathrm{c}}𝑑p_x\left(𝑑p_y𝑑p_z|\psi (\stackrel{}{p},t)|^2\right),$$
(4)
where $`p_\mathrm{c}3p_\mathrm{R}`$ is a good choice since three momentum peaks are initially significantly populated, corresponding to $`2p_\mathrm{R},0,2p_\mathrm{R}`$ pisaBO ; pedri .
Figure 2 shows the initial population in momentum space (inset in (a)) as compared with the population after $`10`$ BO periods, for both the linear and the nonlinear case. The increase of $`C>0`$ has two effects: firstly, it enhances the tunneling for the first few crossings of the BZ. Secondly, it scrambles the out-coupled part of the wave function (see Fig. 2 and its complement in Fig. 4 below), as previously observed in pisaBO ; ingu ; BECBOfirst . The change in the momentum distributions after various Landau-Zener events is a manifestation of the intrinsic instability of the nonlinear GPE dynamics instabil ; wuniu .
Instead of studying the details of the distributions shown in Fig. 2, we will focus on the temporal decay of the survival probability in the following. Figure 3 presents $`P_{\mathrm{sur}}(t)`$, which for the linear case has an exponential form (apart from the $`t0`$ limit raizenN )
$$P_{\mathrm{sur}}(t)e^{\frac{t\mathrm{\Gamma }}{\mathrm{}}},$$
(5)
with the characteristic exponent $`\mathrm{\Gamma }`$. The temporal behavior of $`P_{\mathrm{sur}}`$ depends significantly on $`C`$. For $`C=\pm 0.31`$, we observe clear deviations from a purely exponential decay, as present for small $`C`$. A repulsive nonlinearity initially enhances the tunneling more than after about five crossings of the BZ (see fits to data in Fig. 3). This deviation from the mono-exponential behavior means that the tunneling events occurring at different integer multiples of the Bloch period are correlated by the presence of the nonlinearity. Since the remaining density becomes smaller, the impact of the nonlinearity becomes less. The result is that the rate $`\mathrm{\Gamma }`$ is defined only locally in time, and its value systematically decreases as time increases.
An attractive interaction can stabilize the system at the RET peak, which is shown for $`C=0.31`$ in Fig. 3(b). For optimal comparison, we chose the same initial state (for $`C=+0.31`$) which then was evolved for $`F0`$ with $`C=0.31`$. Such a scenario could be realized by a sudden change of the sign of the scattering length through a Feshbach resonance wieman2001 . This result is consistent with studies of simpler models, where a resonance state can be stabilized at system-specific strengths of the nonlinearity SP2004 ; korsch2005 .
The impact of the nonlinearity on the dynamical evolution of the “closed” system confined to the fundamental BZ can be studied with the help of the recurrence probability Raizen1997 , defined by the autocorrelation
$$P_{\mathrm{rec}}(t)\left|\psi (t)|\psi (t=0)\right|^2.$$
(6)
The BO manifest themselves as the periodic oscillations in $`P_{\mathrm{rec}}(t)`$ plotted in Fig. 4. These oscillations are less and less pronounced with increasing $`C`$, in much the same way as the momentum peaks are washed out when the band edge is crossed in the regime of instability instabil . In contrast to the survival probability, $`P_{\mathrm{rec}}`$ is a phase sensitive measure, and therefore it shows – in addition to the temporal decay – the dephasing of the BO due to the nonlinearity. For $`C=0`$, the recurrence maxima decay in time with the same rate as $`P_{\mathrm{sur}}(t)`$, which offers an alternative method for extracting $`\mathrm{\Gamma }`$. For $`C0`$, $`P_{\mathrm{rec}}`$ can be integrated over time, and the rates are extractable by the approximate proportionality between the integrated area and the inverse decay rate (recalling that $`𝑑tf(t)\mathrm{exp}(t\mathrm{\Gamma })1/\mathrm{\Gamma }`$ to leading order, for a periodic function $`f(t)`$). The latter approach works because we can determine the linear rate from a direct fit to $`P_{\mathrm{rec}}`$ and then compare the ratio of the linear and the nonlinear area (denoted by $`A_0`$ and $`A_C`$). This rough estimate $`\mathrm{\Gamma }_C\mathrm{\Gamma }_0A_0/A_C`$ agrees within $`25\%`$ with the rate extracted from the fits to the data of Fig. 3. The estimate could be improved if we knew the analytic form of the function $`f(t)`$, and it breaks down for large $`C`$, when the periodic oscillations in $`P_{\mathrm{rec}}`$ are destroyed.
Having introduced two methods to extract the tunneling rates, we scan the parameter $`F`$ across a RET peak of the globally exponential curve $`\mathrm{\Gamma }(1/F)`$ (see Fig. 1). The scanned range in $`F`$ corresponds to values of lattice accelerations between $`0.99\mathrm{ms}^2`$ and $`1.65\mathrm{ms}^2`$, which are standard in experiments pisaBO ; instabil .
A repulsive nonlinearity particularly affects the wings of the peak and, for small $`C`$, much less the peak maximum. The global increase of $`\mathrm{\Gamma }`$ with increasing $`C`$ is qualitatively predicted in niu , with enhanced single Landau-Zener crossing probabilities induced by the effective reduction of the energy gap due to the nonlinearity. The left and right-most points in Fig. 1 are in the regime where an amended version of (2) indeed applies niu , and here $`\mathrm{\Gamma }/F`$ is approximately proportional to $`C`$. However, near the peak, the rates do not follow a simple scaling law as a function of $`C`$, and the argumentation of niu does not apply.
For $`C<0`$ we also observe the destruction of the RET peak. For $`C=0.065`$, the BEC clearly stabilizes in the potential wells, whilst for $`C=0.31`$ the situation is more complicated (see Fig. 1). The precise dynamics of the system is governed by the two separate time scales for tunneling and collapse, which strongly depend on parameters in the sensitive RET regime.
In an experiment, $`\omega _x`$ can either be set to zero or decreased to $`\omega _x/2\pi 1\mathrm{Hz}`$ to realize a quasi 1D nonlinear Wannier-Stark problem. We verified that letting $`\omega _x`$ tend to zero for the evolution with $`F0`$, or applying a small finite $`\omega _x`$ gives the same results for the BO cycles studied here. Furthermore, for $`0<C0.05`$, using the renormalized nonlinearity of O1998 we observed that a 1D version of Eq. (3) reproduces well the 3D data. If $`|C|`$ is larger, the nonlinearity couples the longitudinal and transverse degrees of freedom, which affects the dynamics of a real BEC in a non-trivial way Mannella1998 . The 1D computations are feasible up to 100 Bloch periods, and this would allow one to extract the tunneling rates more reliably. The effect of the nonlinearity is, however, hardly visible for $`0<C<0.05`$, and quantitative predictions for a broad range of $`C`$ relied on 3D computations.
To summarize, we observed and quantified the deformation and destruction of the RET peaks due to interactions in a BEC in an accelerated optical lattice. Our results complement ongoing studies of interaction-induced processes such as dynamical instabilities or the decay and subsequent revival of BO. In the regime of small nonlinearity, where dynamical instabilities are not fully developed, the survival and recurrence probabilities experience an exponential decay modified by the condensate nonlinearity. The temporal decay of these observables remains a useful indicator also for large nonlinearity, even if the resonant structure in the tunneling rate is washed out.
We thank M. Cristiani and D. Ciampini for helpful discussions and the Humboldt Foundation (Feodor-Lynen Program), MIUR COFIN-2004, and ESF (QUDEDIS) for support.
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# A phenomenological model of the Resonance peak in High Tc Superconductors
(20/06/2005)
## Abstract
A notable aspect of high-temperature superconductivity in the copper oxides is the unconventional nature of the underlying paired-electron states. The appearance of a resonance peak, observed in inelastic neutron spectroscopy in the superconducting state of the High T<sub>c</sub> cuprates, its apparent linear correlation with the critical superconducting temperature of each of the compounds and its disappearance in the normal state are rather intriguing. It may well be that this peak is the signature of the singlet to triplet excitation, and is an unique characteristic of a d-wave superconductor. We develop a simple criterion for the resonance peak which is based on the concept of twist stiffness and its disappearance at T=T<sub>c</sub>.
The most notable feature of the unconventional nature of the High $`T_c`$ cuprates besides its near-neighbor singlet ground state, is its superconducting gap $`\mathrm{\Delta }`$; unlike conventional $`B.C.S`$ behavior, where $`\mathrm{\Delta }`$, the ampltude of the gap goes up as $`T_c`$ goes up, the measured $`\mathrm{\Delta }`$ as revealed by angular resolved photoemission spectroscopyding $`goes`$ $`down`$ as $`T_c`$ goes up! This had led some authors emery to postulate that the energy scale governing $`T_c`$ is phase stiffness of the order parameter or the superfluid density rather than the modulus of $`\mathrm{\Delta }`$ where $`\mathrm{\Delta }^{}\mathrm{\Delta }`$ is the superconducting condensate density. Inelastic neutron scattering in High T<sub>c</sub> cuprate compounds has been of immense help to enhance our understanding of the magnetic aspects $`underlying`$ physics of High T<sub>c</sub>.regnault It told us right away without any ambiguity that there are at least two clear signatures of the unconventional superconductivity: spin gap and resonance frequency. Some general features emerge from all the compunds so far studied:
(a) Local antiferromagnetic or singlet correlations in the normal and superconducting states are observed , as evidenced by an incoherent background of spin excitation, $`S(q,\omega )`$, particularly at wave vector $`q=\pi ,\pi `$ & frequency $`\omega `$.
(b) On the low energy side , an excitation energy gap , called spin gap $`E_g`$ opens up in the superconducting state, which tends to zero at the critical hole concentration where superconductivity first appears and is a maximum at optimum doping rossatmignod .
(c) The most unexpected feature of the inelastic neutron spectroscopy is the emergence of an extremely intense and narrow peak $`only`$ in the superconducting state at the resonance energy $`\omega _r,`$ at $`q=\pi ,\pi `$ that is a hallmark of each superconducting compound. ($`YBa_2Cu_3O_{7x}`$, siddis $`Bi_2Sr_2CaCu_2O_{8+\delta }`$fong ,$`La_{2x}Sr_xCuO_4`$lake ). This resonance is a collective spin excitation mode where the magnetic excitation spectrum condenses into a peak at a well defined energy. It generally disappears when $`T_c`$ goes to zero and is a generic feature of all the cuprates. The striking characteristic of the resonance peak is its $`linear`$ scaling with $`T_c,`$ as measured for a variety of dopings bourges .
Our main objective in this communication is to convey an underlying universality relating to the resonance peak; the simplicity of the model and its connection to underlying symmetries is its appealing feature.
We assume to start with that the superconducting ground state is a d-wave singlet. In order to bring out the underlying symmetry elements of the superconducting and the normal state, let us introduce the well known concept of superconducting phase stiffness (related to charge stiffness) kohn , spin stiffness shastry as well as that of twist stiffness which is of particular relevance to near neighbor singlets and is associated with chirality. Each of these three stiffnesses are associated with a distinct symmetry operation and expresses the energy increase of the system as each symmetry operation is applied. Let us consider a spinor on site $`i`$
$$\psi _i=\left(\begin{array}{c}c_i^{}\\ c_i^{}\end{array}\right)$$
(1)
Here the $`c_i^{}`$ are the electron creation operator on site $`i`$ in a spin state $``$ and similarly for the other spin $``$. There are three sets of transformation that we can consider on the spinors, one in the charge sector,one in the spin sector and one in the twist sector.
(a) In the charge sector it is given by
$$\psi _i^{^{}}=\mathrm{exp}\left(ie\phi _i\right)\psi _i$$
(2)
where $`e`$ $`is`$ the electron’s charge causing a rotation by an angle $`\phi _i`$, in the electromagnetic gauge space. This is the one parameter transformation of symmetry group $`U(1)`$. In any superconducting ground state, the $`U(1)`$ symmetry will be broken signifying blocking of the phase $`\phi `$ of the superconducting order parameter and hence a non zero superconducting phase stiffness $`D_s.`$ scalapino
(b) We can also rotate the spinor in the spin sector by rotating the spin through an angle $`\theta _i`$ around the spin $`\sigma axis`$ so that
$$\psi _i^{^{}}=\mathrm{exp}\left(i\frac{\sigma }{2}\theta _i\right)\psi _i$$
(3)
where $`\sigma `$ $`is`$ the Pauli spin matrix. The group symmetry is $`SO(3)`$ or $`SU(2)`$. We note that if the ground state is a superconducting $`d`$spin singlet $`S=0,`$ the ground state energy will be unaffected by rotation of the spin axis ,whatever the Hamiltonian is and as a result the spin stiffness $`D_\sigma `$, is necessarily zero .
(c) The twist stiffness is best understood by introducing the chirality where we write
$$\psi _i^{^{}}=\mathrm{exp}\left(\frac{i\sigma \text{ }\gamma _5\theta _i}{2}\right)\psi _i$$
(4)
here the chirality operator $`\gamma _5`$sakurai transcribes the fact that the spin rotation $`\theta _i`$ on site $`i`$ is exactly equal and opposite to that on the near neighbor site $`j`$ whence $`\theta _i\theta _j=2\theta .`$ This gives
$`\psi _i^{^{}}`$ $`=`$ $`\mathrm{exp}\left({\displaystyle \frac{i\sigma \text{ }\theta }{2}}\right)\psi _i\text{ ;}`$ (5)
$`\psi _j^{^{}}`$ $`=`$ $`\mathrm{exp}\left({\displaystyle \frac{i\sigma \text{ }\theta }{2}}\right)\psi _j`$
If the site $`i`$ and $`j`$ belong to sublattice $`A`$ and $`B,`$ then the chiral rotation twists one sublattice around another by a rigid angle $`\theta .`$ The symmetry of the operation because of two sites is $`SU(2)\times SU(2)`$ which is in the same homotopy class as $`SO(4).`$If the ground state is a near neighbor singlet ,the twist rotation $`\theta `$ mixes the singlet with the triplet and hence leads to increase of the ground state energy. This will be clear if we consider the four basis pair states on near neighbor sites $`i`$ & $`j`$ written as sachdev
$``$ $`b={\displaystyle \frac{1}{\sqrt{2}}}(c_i^+c_j^+c_i^+c_j^+)0`$ (6)
$``$ $`t_x={\displaystyle \frac{1}{\sqrt{2}}}(c_i^+c_j^+c_i^+c_j^+)0`$
$``$ $`t_y={\displaystyle \frac{i}{\sqrt{2}}}(c_i^+c_j^++c_i^+c_j^+)0`$
$``$ $`t_z={\displaystyle \frac{1}{\sqrt{2}}}(c_i^+c_j^++c_i^+c_j^+)0`$
Here $`b`$ $`=b^{}0`$ & $`t_\alpha =t_\alpha ^{}0.`$ $`b`$ is a $`S=0`$ singlet while the three $`t_\alpha `$ are $`S=1,`$ triplets and the four states constitute the symmetry $`SO(4)`$. Effect of chiral rotaion on site $`i\&`$ $`j`$ with these basis wave functions. will give $`4\times 4`$ matrix
$$\left(\begin{array}{c}b\\ t_\alpha \end{array}\right)^{^{}}=\mathrm{exp}\left(i\sigma _\alpha \theta \right)\left(\begin{array}{c}b\\ t_\alpha \end{array}\right)$$
(7)
The twist $`\theta `$ has mixed up singlets and triplets. This makes the twist stiffness $`D_t`$ $`0`$ (the twist operation is nonunitary in this ground state). The Lie algebra of $`SO(4)`$ is closed by the three generators that connect $`t_\alpha `$ amongst themselves by pure spin axis rotation,expression (3) and by the other three generators that connect $`b`$ with the three $`t`$ ’s through the twist rotation (5). As for illustration let us a take a model hamiltonian $`tJ`$ hamiltonianzhang
$$H_o=t\underset{i,j}{}c_{i\sigma }^{}c_{j\sigma }J\underset{i,j}{\overset{ij}{}}S_iS_j$$
(8)
Here the first term is the electron hopping between sites $`i`$ and $`j,`$ $`t`$ being the hopping integral while the last term is a Heisenberg antiferromagnetic exchange interaction $`J`$ between first near neighbor spins $`S_i`$ and $`S_j`$ on site $`i`$ and $`j.`$Does $`tJ`$ hamiltonian have a superconducting ground state? The first of the gauge transformation (expression $`2`$ ) has been used to show that its sibling, the Hubbard hamiltonian (in the large $`U`$ limit) has a superconducting ground state with a nonzero phase stiffness $`D_s`$ at $`T=0`$mohit . What about the twist stiffness of the Hamiltonian $`(8)`$? In order to calculate the twist stiffness, we apply uniform twist $`\theta ,`$ between near neighbor sites $`i`$ & $`j`$ of sublattice $`A`$ and $`B`$. It is convenient to transform $`H_o`$ in terms of singlet and triplet pair operators using the pair representation of the spin operatorssachdev which we write as
$`H_o`$ $`=`$ $`t{\displaystyle \underset{i,j}{}}c_{i\sigma }^{}c_{j\sigma }{\displaystyle \frac{3J}{4}}{\displaystyle \underset{\nu }{}}b_\nu ^{}b_\nu +{\displaystyle \frac{J}{4}}{\displaystyle \underset{\nu ,\alpha }{}}t_{\nu \alpha }^{}t_{\nu \alpha }`$ (9)
$`\mu {\displaystyle \underset{i}{}}\left[b_\nu ^2+t_\nu ^2\right]`$
Here $`\mu `$ is the chemical potential assumed same over all space. This term is unaffected by twist and we assume that the sum $`_\nu \left[b_\nu ^2+t_\nu ^2\right]`$ over $`\nu `$ near neighbor pairs which is $`=N_{electron}`$ is conserved. In the limit of small twist the Hamiltonian gets modified
$$H^{^{}}=H\left(0\right)+H\left(\theta \right)$$
(10)
where the first part is the unperturbed untwisted Hamiltonian. By developing $`H\left(\theta \right)`$ to second order we obtain for the perturbing term
$$H\left(\theta \right)=\underset{i,j}{}\left[j_{ij}^\sigma \theta \frac{1}{4}T_{ij}\theta ^2\right]$$
(11)
where $`j_{ij}^\sigma `$ is the spin current operator and $`T_{ij}`$ is the kinetic energy operator. They are given respectively by
$`\widehat{j}`$ $`=`$ $`it[{\displaystyle \underset{ij}{}}(c_{i\sigma }^{}\sigma \gamma _5c_{j\sigma }H.C.)iJ{\displaystyle \underset{\nu }{}}(b_\nu ^{}t_{\nu \alpha }t_{\nu \alpha }^{}b_\nu )]`$
$`\widehat{T}`$ $`=`$ $`t{\displaystyle \underset{ij}{}}(c_{i\sigma }^{}c_{j\sigma }+H.C){\displaystyle \frac{J}{2}}{\displaystyle \underset{\nu }{}}(b_\nu ^2+t_\nu ^2)`$ (12)
We get ground state energy shift due to twist $`\theta `$ as
$$\mathrm{\Delta }E_o=H\left(\theta \right)=\frac{N}{2}D_t\theta ^2$$
(13)
$`D_t\left(\omega =0\right)`$ is the twist stiffness (in two dimensions it has the dimension of energy). It is formally given by
$$D_t\left(\omega \right)=\frac{1}{N}\left[\widehat{T}\underset{n0}{}\left(\frac{0\widehat{j}^\sigma n^2}{ϵ_nϵ_o\mathrm{}\omega }\frac{n\widehat{j}^\sigma 0^2}{ϵ_oϵ_n\mathrm{}\omega }\right)\right]$$
(14)
In the absence of the hopping term and of the spin current term, the energy increase per electron is precisely $`J`$ which is the bare twist stiffness. The first term of $`D_t\left(\omega \right)`$ is the diamagnetic current contribution to stiffness due to the average value of the kinetic energy while the second term reflects second order contribution of $`\mathrm{`}\mathrm{`}`$ paramagnetic spin current conductivity$`{}_{}{}^{\mathrm{"}}\sigma _{t}^{}(\omega )`$ although $`j^\sigma =0`$. The energy levels $`ϵ_n`$ are the triplet excited states for a momentum transfer $`\pi ,\pi `$ (which has a gap $`E_g`$ as measured by inelastic neutron spectroscopy). The spin current in the twisted frame is the response to a “twist vector potential” (engendered by local twist) $`just`$ as the charge current is response to an electromagnetic vector potential. The linear coefficient of the total response is the corresponding twist stiffness. We can rewrite the expression (14) more conveniently in analogy to the missing area sum ruletinkham of the missing Drude weight as
$$D_t\delta (\omega )=D_t^o\delta (\omega )_o^{\mathrm{}}\sigma _t(\omega )𝑑\omega $$
(15)
Here the second term on the right reflects the exhaustion of twist rigidity through incoherent spin excitation where $`\sigma _t(\omega )`$ $`is`$ $`Im\chi _{}\left(\omega \right),`$ the trasverse spin susceptibility. From the experimental neutron data bourges , we know that $`Im\chi _{}\left(\omega \right)is`$ very large at the critical hole concentration $`_h^c`$ at which $`T_c=0`$ while $`Im\chi _{}\left(\omega \right)`$ monotonically decreases ( integrated spectral weight) in the superconducting state as optimum doping $`_h^{opt}`$is approached so that we can reasonably conclude that $`D_t=0`$ at $`_h=_h^c`$ while $`D_t`$ ought to be a maximum at $`_h=_h^{opt}.`$In otherwords $`D_t`$ is a correct indicator of d-wave superconductivity. The non-zero phase stiffness in conventional $`swave`$ superconductor results from broken $`U(1)`$ electromagnetic gauge symmetry. The non-zero spin stiffness in a system with long range magnetic order is associated with a broken $`SO(3)`$ symmetry of the rotational invariance of the spin space and $`D_\sigma `$ goes to zero at $`T=T_N`$ when the invariance is restored. What symmetry or symmetries are broken when the phase coherent singlet d-wave ground state emerges? We may think of the $`dwave`$ superconducting state as a state where $`SO(4)`$ symmetry is explicitly broken as well as $`U(1)`$. The normal state is then a state with zero twist stiffness where the broken $`SO(4)`$ symmetry pertaining to singlet and the three triplets has been restored. If now we accept the premise that at $`T_c`$, twist stiffness $`D_t`$ goes to zero,then one makes the simple statement that $`kT_c`$ is equal to the value of twist stiffness at $`T=0`$ (strictly speaking one should use renormalised stiffness due to triplet excitations) and we have
$$kT_c=D_t\left(T=0,\omega =0\right)$$
(16)
The expression relating spin stiffness to some characteristic frequency (which we shall baptise resonance frequency $`\omega _r`$) can be written as halperin legett
$$D_t\left(T,\omega =0\right)=\chi _{}\left(T\right)\mathrm{}^2\omega _r^2\left(T\right)$$
(17)
The resonance frequency $`\omega _r`$ is a small amplitude harmonic twist oscillation or rigid precession of sublattice $`A`$ with respect to sublattice $`B.`$ Here $`\chi _{}\left(T\right)`$ is the transverse spin flip magnetic susceptibility, which has its largest value at $`Q=\pi ,\pi `$. The transverse static susceptibility $`\chi _{}\left(T\right)`$ in the High $`T_c`$ cuprates (as measured by N.M.R $`\frac{1}{T_{2G}}`$ spin -spin relaxation rate) can be parametrised as julien
$$\chi _{}\left(T\right)=\frac{A}{k\left(T+T_c\right)}$$
(18)
where $`A`$ is a phenomenogical constant. That this form of the static susceptibility in the normal state at $`TT_c`$ is appropriate can be checked from the imaginary part of susceptibility
$$Im\chi _{}(\omega ,T)\frac{\omega }{T}$$
(19)
which is of a form universally observed for small $`\frac{\omega }{T}`$keimer . This behavior in the normal state probably points to proximity to a quantum critical point for spin excitation. In a temperature range above $`T_c`$ the spin correlations have a rapid decay in space but a slow decay in time due to a large density of $`S=1`$ excited states. Real part of the dynamical susceptibility $`\chi _{}(q,\omega )`$ would not show a narrow peak around a specific ordering vector but $`Im\chi _{}(T,\omega )`$ will exhibit considerable weight at low frequency. Using expression (17) and (18) ,we obtain
$`\mathrm{}\omega _r`$ $`=`$ $`akT_c`$ (20)
This is our central result. It corroborates a posteriori the central assumpion that the coherent part of the spectral weight is very simply related to twist stiffness and emerges as coherent resonance peak; and whose intensity reflects the incoherent spin excitations that have disappeared from the energy range $`0\omega E_g`$ in accordance with the sum rule (15). Superconductivity can only arise as the $`dwave`$ singlets manage to shake off the triplets from the normal state soup of singlets& triplets, as a result of the opening up of the spin gap $`E_g(\pi ,\pi )`$. The expression (23) is plotted in figure (1), with neutron and ARPEScampuzano data superimposed. The proportionality constant $`a`$ measured from fig 1 gives the number $`0.42`$ $`me\upsilon /{}_{}{}^{}K`$. If we are at the critical hole doping concentration $`_h^c`$ at which both $`T_N`$,the Neel temperature $`\&`$ $`T_c,`$ the superconducting critical temperature are both zero, then we must have $`D_\sigma =0`$ and $`D_t=0,`$ signifying no long range magnetic order and no long range superfluid order; it is a quantum critical point. It is well known that $`Zn`$ doping destroys $`T_c.`$ It is seen by neutrons that doping with $`Zn`$ introduces large low energy spin fluctuations(integrated spectral weight increaes,the spin gap $`E_g(\pi ,\pi )`$ rapidly goes to zero), that will drive $`D_t`$ to zero suppressing $`\omega _r`$ and killing superconductivity. The normal state can be defined as a spin liquid (by definition has no sublattice magnetisation) where we have considerable low energy spin excitation. We also require that translational invariance be unbroken for the system to qualify as a liquid. Thus it describes a gapless spin liquid more in conformity with the original suggestion of the long range RVB liquid anderson In its loss of twist stiffness the spin liquid behaves like any conventional liquid loosing shear rigidity at the melting transition. The concept of twist stiffness is based on $`infinitesimally`$ $`small`$ $`twist`$ as is customary in these definitions; beyond $`TT_c,`$ the restored dynamical $`SO(4)`$ symmetry implies $`bt_\alpha `$ pair fluctuation in the spin liquid phase costing no energy around the untwisted singlet. If this symmetry persists for $`all`$ twist angles then we will be in the frustrated “Henley limit”henley of infinite classical spin degeneracy where one sublattice $`A`$ will twist freely around the other sublattice $`B`$ and the two sublattices are totally decoupled. Or else the system may develop a region where $`D_t\left(TT_c\right)`$ $`may`$ $`become`$ $`negative`$ for $`large`$ $`twist`$ $`angles`$ generating large singlet-triplet excursions and hence may go spontaneously to a distorted or twisted ground state siggia . Although twist stiffness and superconducting phase stiffness are different at $`T=0,`$ their simultaneous disappearance at $`T=T_c`$ is indicated by the Arpes results campuzano of the hump and dip structure in the electronic spectral weight and point to strong coupling of triplet and phase fluctuation as $`T_c`$ is approached.
Several theoretical models exist maki that explain the resonance peak. Our objective in this paper has been relatively simple: can we understand the resonance peak without a detailed model and does it have some predictive ability as to the underlying symmetry nature of the normal and superconducting state? I think the arguments given in this paper will throw some new light on these issues.
The author acknowledges highly stimulating discussions with Dr M. Cuoco, He benefited from comments by M. Avignon, P. Bourges, T. Chatterjee, K. Jain, J. Ranninger, T. Ziman, The author is very grateful to Dr S. Fratini for a careful reading of the manuscript at various stages.
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# Hitting time and dimension in Axiom A systems and generic interval excanges
## 1 Introduction
It is well known by classical recurrence results that a typical orbit of a dynamical systems comes back (in any reasonable neighborhood) near to its starting point. The quantitative study of recurrence quantifies the speed of this coming back, estimating, for example, how much time is needed to come back in any ball centered in the starting point (the reader can find and exposition of more and less recent developments about this kind of questions in the survey \[CG\]). It turns out that in many cases the scaling law of return times are related to the dimension of the invariant measure of the system. More precisely, let us consider a starting point $`x`$, a ball $`B(x,r)`$ and the time $`\tau _{B(x,r)}(x)`$ needed for the starting point $`x`$ to come back to $`B(x,r).`$ With these notations we have, for example, (\[S\], \[STV\]) that in exponentially mixing systems or positive entropy (with some small tecnical assumptions) systems over the interval $`\tau _{B(x,r)}(x)r^{d(x)}`$, where $`d(x)`$ is the local dimension at $`x.`$ Moreover (\[Bo\], \[BS\]) in general finite measure preserving systems the recurrence gives a lower bound to the dimension ($`\underset{r0}{lim}\frac{\mathrm{log}\tau _{B(x,r)}(x)}{\mathrm{log}r}d(x)`$).
A similar and strictly related (see e.g. \[LHV\]) problem is about the time needed for a typical point $`x`$ of an ergodic system to enter in some neighborhood of another point $`y.`$This leads to the *hitting time* (also called *waiting time*)indicators. For example let us denote by $`\tau _{B(y,r)}(x)`$ the time needed for the point $`x`$ to enter in the ball $`B(y,r)`$ with center $`y`$ and radius $`r`$ (this in some sense *generalizes* the recurrence because we are allowed to consider an arriving point $`y`$ different from the starting point $`x`$). Again we consider the scaling behavior of this waiting time for small $`r`$. The waiting time indicator will have value $`R`$ if $`\tau _{B(y,r)}(x)r^R`$ (precise definitions in section 2). Again, there are relations with the local dimension. Some general relations are proved in \[G\] (see thm. 3) and show that the waiting time indicator gives an upper bound to the dimension. Moreover there is a class of systems where the waiting time is equal to dimension. This class of systems includes for example systems having exponential distribution of return times in small balls (this includes many *more or less* hyperbolic systems over the interval \[BSTV\]). We remark that exponential return times in balls is conjecured but yet not proved in general Axiom A systems, thus equality between hitting time and dimension for axiom A does not follow from such result.
We want to remark that there are also some relevant cases where the equality between recurrence or hitting time with dimension does not hold, hence this kind of questions are not trivial. Such cases includes rotations by Liuouville numbers (see \[BS\],\[KS\]), and Maps having an indifferent fixed point and infinite invariant measure (\[GKP\]).
A further motivation for this kind of studies is that the relations between recurrence (and similar) with dimension are used in the physical literature \[HJ\],\[GE\],\[JKLPS\] to provide numerical methods for the study of the Hausdorff dimension of attractors. Since recurrence gives a lower bound to dimension and hitting time gives an upper bound, the combined use of these may produce efficient numerical estimators for the dimension of attractors.
The main result of this note is to show that in nontrivial nice examples such as Axiom A systems and typical interval exchanges, the hitting time indicator equals the local dimension $`d(y)`$ of the considered measure. Hence, in such systems, for typical $`x`$ and $`y`$ we will have $`\tau _{B(y,r)}(x)r^{d(y)}`$. As an application of these results we give an estimation for the Birkoff averages of functions (in ergodic systems) having some asymptote and no finite $`L^1`$norm. Here the Birkoff average will increase to infinity as the number of iterations increases (this is trivially by ergodic theorem). The hitting time indicator will give an estimation about the speed of going to infinity of such average.
###### Acknowledgement 1
I wish to thank Corinna Ulcigrai , Jean Rene Chazottes and Stefano Marmi for fruitful discussions, which allowed me to discover some relevant literature and to simplifiyng much the proof of the main result.
## 2 Generalities and a criteria for hitting time and dimension
In the following we will consider a discrete time dynamical system $`(X,T)`$ were $`X`$ is a separable metric space equipped with a Borel finite measure $`\mu `$ and $`T:XX`$ is a measurable map.
Let us consider the first entrance time of the orbit of $`x`$ in the ball $`B(y,r)`$ with center $`y`$ and radius $`r`$
$$\tau _r(x,y)=\mathrm{min}(\{n𝐍,n>0,T^n(x)B(y,r)\}).$$
By considering the power law behavior of $`\tau _r(x,y)`$ as $`r0`$ let us define the hitting time indicators as
$$\overline{R}(x,y)=\underset{r0}{limsup}\frac{\mathrm{log}(\tau _r(x,y))}{\mathrm{log}(r)},\underset{¯}{R}(x,y)=\underset{r0}{liminf}\frac{log(\tau _r(x,y))}{\mathrm{log}(r)}.$$
If for some $`r,`$ $`\tau _r(x,y)`$ is infinite then $`\overline{R}(x,y)`$ and $`\underset{¯}{R}(x,y)`$ are set to be equal to infinity. The indicators $`\overline{R}(x)`$ and $`\underset{¯}{R}(x)`$ of quantitative recurrence defined in \[BS\] are obtained as a special case, $`\overline{R}(x)=\overline{R}(x,x)`$, $`\underset{¯}{R}(x)=\underset{¯}{R}(x,x)`$.
We recall some basic properties of $`\overline{R}(x,y):`$
###### Proposition 2
$`R(x,y)`$ satisfies the following properties
* $`\overline{R}(x,y)=\overline{R}(T(x),y)`$, $`\underset{¯}{R}(x,y)=\underset{¯}{R}(T(x),y)`$.
* If $`T`$ is $`\alpha Hoelder`$, then $`\overline{R}(x,y)\alpha \overline{R}(x,T(y))`$, $`\underset{¯}{R}(x,y)\alpha \underset{¯}{R}(x,T(y))`$.
* If we consider $`T^n`$ instead of $`T`$ $`\overline{R}_T(x,y)\overline{R}_{T^n}(x,y)`$, $`\underset{¯}{R}_T(x,y)\underset{¯}{R}_{T^n}(x,y)`$.
*Proof.* The first two points comes from \[G\] (and they comes directly from definitions). For the third one let us denote with $`\tau `$ and $`\tau ^{}`$the hitting time with resp. to $`T`$ and $`T^n`$. By definition $`\overline{R}_T(x,y)=\underset{r0}{limsup}\frac{\mathrm{log}(\tau _r(x,y))}{\mathrm{log}(r)}`$ but $`\tau _r(x,y)n\tau _r^{}(x,y)`$ and $`\frac{\mathrm{log}(\tau _r(x,y))}{\mathrm{log}(r)}\frac{\mathrm{log}(\tau _r^{}(x,y))+\mathrm{log}n}{\mathrm{log}(r)},`$ and taking the $`limsup`$ we are done. The same can be done for the $`liminf`$.$`\mathrm{}`$
In general systems the quantitative recurrence indicator gives only a *lower* bound on the dimension (\[BS\], \[Bo\]). The waiting time indicator instead give an *upper* bound (\[G\]) to the local dimension of the measure at the point $`y`$. This is summarized in the following
###### Theorem 3
If $`(X,T,\mu )`$ is a dynamical system over a separable metric space, with an invariant measure $`\mu .`$ For each $`y`$
$$\underset{¯}{R}(x,y)\underset{¯}{d}_\mu (y),\overline{R}(x,y)\overline{d}_\mu (y)$$
(1)
holds for $`\mu `$ almost each $`x`$. Where $`\underset{¯}{d}_\mu (y)`$ and $`\overline{d}_\mu (y)`$ are the lower and upper local dimension at $`y.`$<sup>3</sup><sup>3</sup>3If $`X`$ is a metric space and $`\mu `$ is a measure on $`X`$ the upper local dimension at $`xX`$ is defined as $`\overline{d}_\mu (x)=\underset{r0}{limsup}\frac{log(\mu (B(x,r)))}{log(r)}=\underset{k𝐍,k\mathrm{}}{limsup}\frac{log(\mu (B(x,2^k)))}{k}`$ The lower local dimension $`\underset{¯}{d}_\mu (x)`$ is defined in an analogous way by replacing $`limsup`$ with $`liminf`$. If $`\overline{d}_\mu (x)=\underset{¯}{d}_\mu (x)=d`$ almost everywhere the system is called exact dimensional. In this case many notions of dimension of a measure will coincide (see for example the book \[P\]).Moreover, if $`X`$ is a closed subset of $`^𝐧`$, then for almost each $`xX`$
$$\overline{R}(x,x)\overline{d}_\mu (x),\underset{¯}{R}(x,x)\underset{¯}{d}_\mu (x).$$
A natural question which is important also from the numerical applications is whether equality can replace the above inequalities (see the results from \[S\], \[G\], \[BS\], \[STV\] already outlined in the introduction). The following is a general criteria that assures (together with theorem 3) for typical points, equality between waiting time and local dimension.
###### Lemma 4
Let $`xX`$ and
$$SF_r^n(x)=XB(x,r)T^1(XB(x,r))T^2(XB(x,r))\mathrm{}T^n(XB(x,r))$$
be the set of points that after $`n`$ steps never enters into $`B(x,r).`$ If for each $`ϵ>0`$ we have $`\mu (SF_{2^n}^{\mu (B(x,2^n))^{1ϵ}})<\mathrm{}`$ then for almost each $`y`$ it holds $`\overline{R}(y,x)\overline{d}_\mu (x)`$ and $`R`$$`(y,x)\underset{¯}{d}_\mu (x).`$
*Proof.* The proof follows by a Borel Cantelli argument. Let $`R_ϵ=\{yX,\overline{R}(y,x)(1+ϵ)\overline{d}_\mu (x)\}`$. If we prove that this set has measure zero for each $`ϵ`$ we are done. If we know that for some $`ϵ`$ $`\mu (SF_{2^n}^{\mu (B(x,2^n))^{1ϵ}})<\mathrm{}`$ this means that the set of points such that $`\tau _{2^n}(y,x)>\mu (B(x,2^n))^{1ϵ}`$ for infinitely many $`n`$ has zero measure. Taking logarithms and dividing by $`n`$ we have $`\frac{\mathrm{log}(\tau _{2^n}(x,y))}{n}(1+ϵ)\frac{\mathrm{log}(\mu (B(x,2^n)))}{n}`$ eventually (as $`n`$ increases) for a full measure set and then $`\overline{R}(y,x)=limsup\frac{\mathrm{log}(\tau _{2^n}(x,y))}{n}(1+ϵ)limsup\frac{\mathrm{log}(\mu (B(x,2^n)))}{n}=(1+ϵ)\overline{d}_\mu (x)`$ on a full measure set. This is true for each $`ϵ`$ and we have the statement. The same can be done for the proof of $`R`$$`(y,x)\underset{¯}{d}_\mu (x)\mathrm{}`$
## 3 Axiom A systems
In this section we will consider Axiom A systems, we will apply the properties of Gibbs measures to prove that they satisfy Lemma 4 at almost all points. We will prove the following
###### Theorem 5
If $`X`$ is a basic set of an axiom A diffeomorphism, $`\mu `$ is an equilibrium measure for an Hoelder potential defined on $`X`$. Then $`(X,T,\mu )`$ satisfies Lemma 4 at almost each $`xX`$ and hence for almost each $`x`$ it holds $`R`$$`(y,x)=\underset{¯}{d}_\mu (x),\overline{R}(y,x)=\overline{d}_\mu (x)`$ for almost each $`y.`$
First we need a general estimation on the behavior of a certain kind of sequences.
###### Lemma 6
Let $`0<m<1,`$and $`a_n`$ be defined by $`\{\begin{array}{c}a_n=a_{n1}m+s_n\\ a_0=m^2\end{array}`$<sup>2</sup><sup>2</sup>footnotetext: Keywords: dimension, quantitative recurrence, Axiom A, interval exchanges, Birkoff sums where $`s_n=\frac{2n+1}{n^2(n+1)^2}=\frac{1}{i^2}\frac{1}{(i+1)^2}`$ then for $`n2`$ it holds $`a_n\frac{m^{[\frac{n}{2}]}}{1m}+\frac{4}{n^2}.`$
*Proof.* We have
$$a_n=m^{n+1}+m^{n1}s_1+m^{n2}s_2+m^{n3}s_3+\mathrm{}+ms_{n1}+s_n.$$
Since $`s_i<1`$ and $`m<1`$ then $`a_n_{n/2}^nm^i+_{n/2}^ns_i=\frac{m^{[\frac{n}{2}]}m^n}{1m}+\frac{1}{([n/2])^2}\frac{1}{([n/2]+1)^2}\frac{m^{[\frac{n}{2}]}}{1m}+\frac{4}{n^2}.\mathrm{}`$
*Proof of Theorem 5.* We already know that (thm. 3) $`\underset{¯}{R}(x,y)\underset{¯}{d}_\mu (y),\overline{R}(x,y)\overline{d}_\mu (y)`$. For the opposite inequalities, first we remark that (see \[Bow\] pag. 72) $`X=X_1\mathrm{}X_l`$ where $`T(X_i)=X_{i+1},(T(X_l)=X_1)`$ and $`T^l|_{X_i}`$ is topologically mixing.
By Lemma 2 we can suppose that $`x,y`$ belongs to the same $`X_i.`$ and that $`T`$ is topologically mixing (replacing $`T`$ with $`T^l`$ we have a mixing transformation on $`X_i`$, moreover by the fourth point of Lemma 2 we see that if we have an upper bound for $`\overline{R}_{T^l}`$ and $`R`$$`_{T^l}`$then this is also an upper bound for $`\overline{R}_T(x,y)`$ and $`R`$$`{}_{T}{}^{}(x,y).`$ Since in this proof we are looking for an upper bound, by replacing $`T`$ with $`T^l`$ we can suppose that the map is topologically mixing).
To estimate the measure of the set $`SF_r^n(x)`$ let us consider a Markov partition $`Z=\{Z_i\}`$ of $`X.`$ Let $`Z_m^n=T^m(Z)\mathrm{}T^n(Z).`$ By uniform hyperbolicity there are constants $`C,\lambda >0`$ such that $`diam(Z_n^n)Ce^{\lambda n}.`$ By this we know that there is some
$$n(r)\lambda ^1(\mathrm{log}r\mathrm{log}C)$$
such that the partition is of size so small that there is one element $`Z_0`$ of the partition $`\overline{Z}=Z_{n(r)}^{n(r)}`$ which is included in $`B(x,r).`$
Now $`SF_r^n(x)B_0^n=XZ_0T^1(XZ_0)\mathrm{}T^n(XZ_0).`$ We remark that $`B_0^n`$ is the union of many cylinders. The measure of $`B_0^n`$ decreases very fast by the weak Bernoulli property of the equilibrium measure $`\mu .`$ Indeed by \[Bow\] pag. 90, we know that since the map $`T`$ can be supposed to be topologically mixing and then $`\mu `$ has the weak Bernoulli property: i.e. let us consider $`t,s0`$ and the partitions $`P_s=ZT^1(Z)\mathrm{}T^s(Z)`$ and $`Q_t=T^t(Z)\mathrm{}T^{tk}(Z).`$ For each $`ϵ,`$ if $`ts=N_Z(ϵ)`$ is big enough, then $`\underset{PP_s,QQ_t}{}\mu (QP)\mu (P)\mu (Q)<ϵ.`$ Moreover by \[Bow\], theorem 1.25 we can find an estimation for $`N_Z(ϵ)`$ as a function of $`ϵ`$ (see \[Bow\] pag. 38): $`N_Z(ϵ)=c\mathrm{log}(ϵ)+c^{},`$ where $`c,c^{}`$ are constants depending on $`\mu ,T`$ and $`Z`$.
The estimation for $`\mu (B_0^l)`$ follows from the fact that a non empty cylinder for the partition $`\overline{Z}=Z_{n(r)}^{n(r)}`$ is also a cylinder for the partition $`Z.`$
Indeed the cylinder $`\overline{𝐳}_m=\overline{Z}_{i_1}T^1(\overline{Z}_{i_2})\mathrm{}T^{m1}(\overline{Z}_{i_m})`$ , $`\overline{Z_i}\overline{Z}`$ satisfies $`\overline{𝐳}_m=𝐳_{m+2n}`$ where $`𝐳_{m+2n}=T^n(Z_{j_1})T^{n1}(Z_{j_2})\mathrm{}T^{m1n}(Z_{j_{(m+2n)}})`$ is a cylinder of $`Z`$ and $`\overline{Z}_{i_k}=T^{k+n}Z_{j_k}\mathrm{}T^{kn}Z_{j_{k+2n}}.`$ Since $`\mu `$ is preserved then $`\mu (T^{k+n}Z_{j_k}\mathrm{}T^{kn}Z_{j_{k+2n}})=\mu (T^kZ_{j_k}\mathrm{}T^{k2n}Z_{j_{k+2n}})`$ hence we can apply the weak Bernoulli property to such cylinders obtaining that also $`\overline{Z}`$ satisfies such a property and
$$N_{\overline{Z}}(ϵ)c\mathrm{log}(ϵ)+c^{}+2n$$
(2)
(where $`c,c^{}`$ are the constants of $`N_Z(ϵ)`$ as above). We recall that $`n`$ depends on $`r`$ and we can choose $`n(r)\lambda ^1(\mathrm{log}r\mathrm{log}C).`$
Finally let us apply the weak Bernoulli property of $`\overline{Z}`$ to get an estimation for $`\mu (B_0^n).`$ Let us set $`m=\mu (XZ_0)`$ and $`ϵ(i)=\frac{2i+1}{i^2(i+1)^2}=\frac{1}{i^2}\frac{1}{(i+1)^2}`$. We have (eq. 2) that setting $`C^{}(r)=c^{}2\lambda ^1(\mathrm{log}r\mathrm{log}C)`$ then $`N_{\overline{Z}}(ϵ)c\mathrm{log}(ϵ)+c^{}+2n=c\mathrm{log}(\frac{2i+1}{i^2(i+1)^2})+C^{}(r)`$ and there is a $`C`$ s.t. $`N_{\overline{Z}}(ϵ)C\mathrm{log}(i)+C^{}(r).`$
Let us set $`n_i=_{ji}N_{\overline{Z}}(ϵ(j)).`$ Thus for each $`\delta `$ there is a $`K`$ such that, if $`i`$ is big enough
$$n_iKi^{1+\delta }\mathrm{log}r.$$
(3)
The measure of $`B_0^{n_i}`$ can then be estimated applying $`i`$ times the Bernoulli property, with $`ϵ(i)=`$ $`\frac{2i+1}{i^2(i+1)^2}`$ as above, to subcylinders of increasing length $`N_{\overline{Z}}(ϵ(1)),N_{\overline{Z}}(ϵ(1))+N_{\overline{Z}}(ϵ(2)),N_{\overline{Z}}(ϵ(1))+N_{\overline{Z}}(ϵ(2))+N_{\overline{Z}}(ϵ(3))\mathrm{}`$ obtaining by the Bernoulli property of $`\mu `$
$`\mu (B_0^{N(ϵ(1))})`$ $``$ $`m^2+ϵ(1),`$
$`\mu (B^{N(ϵ(1))+N(ϵ(2))})`$ $``$ $`(m^2+ϵ(1))m+ϵ(2),`$
$`\mu (B^{N(ϵ(1))+N(ϵ(2))+N(ϵ(3))})`$ $``$ $`((m^2+ϵ(1))m+ϵ(2))m+ϵ(3),\mathrm{}`$
Hence by Lemma 6 above
$$\mu (B_0^{n_i})\frac{m^{[\frac{i}{2}]}}{1m}+\frac{4}{i^2}.$$
We remarked that $`SF_r^nB_0^n.`$ If we consider another element $`Z_1`$ of $`\overline{Z}`$ with $`Z_1B(x,r)`$ and $`B_1^n=`$ $`XT^1(X(Z_0Z_1))\mathrm{}T^n(XZ_0Z_1)),`$ we have also $`SF_r^nB_1^iB_0^l.`$ Now considering a sequence $`Z_0,\mathrm{},Z_w`$ of elements of $`\overline{Z}`$ with $`Z_0,\mathrm{},Z_wB(x,r)`$ and $`B_w^n=XT^1(X(Z_0\mathrm{}Z_w))\mathrm{}T^n(X(Z_0\mathrm{}Z_w)),`$ we have also $`SF_r^nB_w^n.`$ The measure of $`B_w^n`$ can be estimated as above, obtaining $`\mu (B_w^{n_i})\frac{m_w^{i/2}}{1m_w}+\frac{4}{i^2},`$ where $`m_w=\mu (X(Z_0\mathrm{}Z_w)).`$
Now, refining again the partition $`\overline{Z}`$ if necessary (this is true because the diameter of each $`Z_i`$ is less or equal than $`Ce^{\lambda n},`$ and this will only change the constants in $`N(ϵ)`$) we can suppose that the diameter of each piece of the partition has diameter less than $`r/4.`$ We then have that we can choose $`Z_0,\mathrm{},Z_w`$ such that $`B(x,\frac{r}{2})Z_0\mathrm{}Z_wB(x,r).`$ Then $`\mu (X(Z_0\mathrm{}Z_w))\mu (XB(x,\frac{r}{2})).`$This gives,
$$\mu (B_w^{n_i})\frac{(1\mu (B(x,\frac{r}{2})))^{\frac{i}{2}}}{\mu (B(x,\frac{r}{2}))}+\frac{4}{i^2}.$$
By 3 we then have $`i\frac{n_i^{\frac{1}{1+\delta }}}{(K\mathrm{log}r/4)^{^{\frac{1}{1+\delta }}}}`$, by this
$`\mu (SF_{2^n}^{\mu (B(x,2^n))^{1ϵ}})`$ $``$ $`\mu (B_w^{\mu (B(x,2^n))^{1ϵ}})`$
$``$ $`{\displaystyle \frac{(1\mu (B(x,2^{n1})))^{(Kn+\mathrm{log}4)^{\frac{1}{1+\delta }}\mu (B(x,2^n))^{\frac{1ϵ}{1+\delta }}}}{\mu (B(x,2^{n1}))}}+`$
$`+{\displaystyle \frac{4}{(Kn+\mathrm{log}4)^{\frac{2}{1+\delta }}\mu (B(x,2^{n1}))^{^{\frac{22ϵ}{1+\delta }}}}}.`$
Since in our case $`d_\mu (x)=d<\mathrm{}a.e.,`$ there is a constant $`Q,`$ (which depends on the local dimension) such that $`0<Q<\frac{\mu (B(x,2^{n1}))}{\mu (B(x,2^n))}<1`$ when $`n`$ is big, recalling that $`\delta `$ can be chosen as small as we want and hence smaller than $`ϵ`$, then $`\mu (SF_{2^n}^{\mu (B(x,2^n))^{1ϵ}})`$ is less than about $`\frac{(e^{\frac{Q}{2}})^{Kn^{\frac{1}{1+\delta }}\mu (B(x,2^n))^{\frac{ϵ\delta }{1+\delta }}}}{\mu (B(x,2^{n1}))}+4(Kn+\mathrm{log}4)^{\frac{2}{1+\delta }}\mu (B(x,2^{n1}))^{^{\frac{2+2ϵ}{1+\delta }}}`$ and we have $`\underset{𝑛}{}\mu (SF_{2^n}^{\mu (B(x,2^n))^{1ϵ}})<\mathrm{}.`$This is enough to apply the Lemma 4 and have the required statement.$`\mathrm{}`$
## 4 Interval exchanges
An interval excange is a piecewise isometry which preserves the Lesbegue measure. In this section we apply a result of Boshernitzan about a full measure class of uniquely ergodic interval exchanges maps to prove equality between hitting time and dimension at discontinuity points. We refer to \[Bo2\] for generalities on this important class of maps.
###### Theorem 7
For a typical interval exchange transformation $`T`$ (for a full measure set, in the space of interval exchanges) for each discontinuity point $`x_0`$ it holds $`R`$$`(y,x_0)=1`$ for almost each $`y[0,1].`$
*Proof.* By a result of (\[Bo2\]) we have that if $`T`$ is a typical i.e.t. and $`\delta (n)`$ is the minimum distance between the discontinuity points of $`T^n,`$ then there is a constant $`C`$ and a sequence $`n_k`$ such that $`\delta (n_k)\frac{C}{n_k}.`$ If $`x_0`$ is a discontinuity point then for each $`n_k`$ it also hold that $`\mathrm{min}_{i,jn_k}d(T^i(x_0),T^j(x_0))\frac{C}{n_k}.`$ Let us consider the set $`J_k=_{in_k}B(T^i(x_0),\frac{C}{3n_k})`$. Since it is a union of disjoint balls the measure of $`J_k`$ is $`\frac{2C}{3}.`$ This implies that in the interval $`[0,1]`$ there is a positive measure set $`J`$ of points belonging to infinitely many $`J_k.`$
If a point $`y`$ is in $`J`$ then for a subsequence $`n_{k_i}`$ it holds $`d(T^{n_{k_i}}(y),x_0)\frac{C}{3n_{k_i}}`$. This is true because by the Boshernitzan result there are no counter images of other discontinuity points in the interval $`[y,T^n(x_0)]`$ and then these two points cannot be separated during $`n`$ iterations of the map.
Since $`liminfnd(T^n(y),x_0)\frac{C}{3}`$ then $`R`$$`(y,x_0)1`$ and we have the statement for $`y`$ varying in a positive measure set $`J`$. Since the system is ergodic, by Proposition 2 we have the statement.$`\mathrm{}`$
Since in interval exchanges th only source of initial condition sensitivity is the discontinuity (the orbits of two points can be only separated by a discontinuity) we remark that an estimation of the approaching speed of typical orbits to the discontinuity is useful to estimate the kind of ”weak” chaos that is present is such maps. The theorem above in some sense can give (using the construcion done in \[BGI\] ) an upper bound on the initial condition sensitivity of such maps. We will not go into details about this in this work however.
## 5 Hitting time and Birkoff sums
Let us consider a function $`f:X\{x_0\}`$ , $`f0`$ which is continuos and which satisfies $`_Xf𝑑\mu =\mathrm{}`$ because it has an asymptote in $`x_0`$ where $`f(x)d(x,x_0)^\alpha .`$
By the ergodic theorem we know that for almost each $`x`$ the Birkoff average $`\frac{S_n(x)}{n}=\frac{1}{n}_{i=0}^nf(T^i(x))`$ is such that $`\frac{S_n(x)}{n}\mathrm{}.`$ By the results of the previous sections, if we know the hitting time indicator at $`x_0`$ we can have an estimation for the speed of increasing of $`\frac{S_n(x)}{n}.`$
###### Theorem 8
Let us suppose that near $`x_0`$ we have $`0<lim\frac{f(x)}{d(x,x_0)^\alpha }<\mathrm{}`$, for $`\alpha >1`$ then for almost each $`x`$
$$\frac{\alpha }{\overline{R}(x,x_0)}\underset{n\mathrm{}}{limsup}\frac{\mathrm{log}(S_n(x))}{\mathrm{log}(n)}\frac{\alpha }{\underset{¯}{R}(x,x_0)}+1$$
*Proof.* By the definition of $`R`$$`(y,x_0)`$ we obtain (see \[G\]) that for each $`ϵ>0`$
$$\underset{n\mathrm{}}{liminf}n^{\frac{1}{\underset{¯}{R}(y,x_0)}+ϵ}d(T^n(y),x_0)=\mathrm{}.$$
Then we have that if $`n`$ is big enough $`d(T^n(y),x_0)n^{\frac{1}{\underset{¯}{R}(y,x_0)}ϵ}.`$ Now we remark that since $`X`$ is compact there are $`c_1,c_2`$ such that $`f(x)\mathrm{max}(c_1,c_2d(x,x_0)^\alpha ).`$ Then if $`n`$ is big enough
$`{\displaystyle \underset{i=0}{\overset{n}{}}}f(T^i(y))`$ $``$ $`{\displaystyle \underset{i=0}{\overset{n}{}}}\mathrm{max}(c_1,c_2d(T^i(y)x_0)^\alpha )`$
$``$ $`{\displaystyle \underset{i=0}{\overset{n}{}}}\mathrm{max}(c_1,c_2n^{\frac{\alpha }{\underset{¯}{R}(y,x_0)}+\alpha ϵ})nc_1+c_2n^{\frac{\alpha }{\underset{¯}{R}(y,x_0)}+\alpha ϵ+1}`$
and we have $`\underset{n\mathrm{}}{limsup}\frac{\mathrm{log}(S_n(x))}{\mathrm{log}(n)}\frac{\alpha }{\underset{¯}{R}(y,x_0)}+1.`$ On the other hand, by the definition of $`\overline{R}(y,x_0)`$ we have that frequently $`d(T^n(y),x_0)n^{\frac{1}{\overline{R}(y,x_0)}+ϵ}`$, then frequently $`_{i=0}^nf(T^i(y))cn^{\frac{\alpha }{\overline{R}(y,x_0)}\alpha ϵ}`$.$`\mathrm{}`$
By the above result and the previous ones it easily follows :
1. (by thm. 3) In a general system, if the local dimension at $`x_0`$ is $`d_\mu (x_0).`$ Then for almost each $`x`$
$$\underset{n\mathrm{}}{limsup}\frac{\mathrm{log}(S_n(x))}{\mathrm{log}(n)}\frac{\alpha }{d_\mu (x_0)}+1$$
2. (by thm. 7) If $`T`$ is an IET and $`x_0`$ is a discontinuity point then
$$\alpha \underset{n\mathrm{}}{limsup}\frac{\mathrm{log}(S_n(x))}{\mathrm{log}(n)}\alpha +1$$
3. (by thm. 5) If $`(X,T)`$ is axiom A (with an equilibrium measure, as in thm. 5), $`x_0`$, $`x`$ are typical and $`d`$ is the dimension of the measure then
$$\frac{\alpha }{d}\underset{n\mathrm{}}{limsup}\frac{\mathrm{log}(S_n(x))}{\mathrm{log}(n)}\frac{\alpha }{d}+1.$$
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# EGRET Excess of Galactic Gamma Rays as Signal of Dark Matter Annihilation
## 1 Introduction
Cold Dark Matter (CDM) makes up 23% of the energy of the universe, as deduced from the WMAP measurements of the temperature anisotropies in the Cosmic Microwave Background, in combination with data on the Hubble expansion and the density fluctuations in the universe $`^\mathrm{?}`$. The Dark Matter has to be much more widely distributed than the visible matter, since the rotation speeds do not fall off like $`1/\sqrt{r}`$, as expected from the visible matter in the centre, but stay more or less constant as function of distance. For a ”flat” rotation curve the DM has to fall off slowly like $`1/r^2`$ instead of the exponential drop-off for the visible matter. The fact that the DM is distributed over large distances implies that its properties must be quite different from the visible matter, since the latter clumps in the centre owing to its rapid loss of kinetic energy by the electromagnetic and strong interactions after infall into the centre. Since the DM apparently undergoes little energy loss, it can have at most weak interactions. In addition its mass is probably large, since it cannot be produced with present accelerators. Therefore it is generically called a WIMP, a Weakly Interacting Massive Particle.
Weakly interacting particles can annihilate, yielding predominantly quark-antiquark pairs in the final state, which hadronize into mesons and baryons. The stable decay and fragmentation products are neutrinos, photons, protons, antiprotons, electrons and positrons. From these, the protons and electrons disappear in the sea of many matter particles in the universe, but the photons and antimatter particles may be detectable above the background, generated by particle interactions. Such searches for indirect Dark Matter detection have been actively pursued, see e.g the review by Bergstr$`\ddot{\mathrm{o}}`$m $`^\mathrm{?}`$ or more recently by Bertone, Hooper and Silk.$`^\mathrm{?}`$
The present analysis on diffuse galactic gamma rays differs from previous ones by considering simultaneously the complete sky map and the energy spectrum, which allows us to constrain both the halo distribution and the WIMP mass. More details have been given elsewhere $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$. The constraint on the WIMP annihilation cross section from WMAP is discussed in Section 2, while the constraints on the mass and the DM halo profile from the EGRET excess are discussed in Sections 3. The summary is given in Section 4.
## 2 Annihilation Cross section Constraints from WMAP
In the early universe all particles were produced abundantly and were in thermal equilibrium through annihilation and production processes. At temperatures below the mass of the WIMPS the number density drops exponentially. The annihilation rate $`\mathrm{\Gamma }=<\sigma v>n_\chi `$ drops exponentially as well, and if it drops below the expansion rate, the WIMP’s cease to annihilate. They fall out of equilibrium (freeze-out) at a temperature of about $`m_\chi /22`$ $`^\mathrm{?}`$ and a relic cosmic abundance remains.
For the case that $`<\sigma v>`$ is energy independent, which is a good approximation in case there is no coannihilation, the present mass density in units of the critical density is given by $`^\mathrm{?}`$:
$$\mathrm{\Omega }_\chi h^2=\frac{m_\chi n_\chi }{\rho _c}(\frac{210^{27}cm^3s^1}{<\sigma v>}).$$
(1)
One observes that the present relic density is inversely proportional to the annihilation cross section at the time of freeze out, a result independent of the WIMP mass (except for logarithmic corrections). For the present value of $`\mathrm{\Omega }_\chi h^2=0.113\pm 0.009`$ the thermally averaged total cross section at the freeze-out temperature of $`m_\chi /22`$ must have been around $`210^{26}\mathrm{cm}^3\mathrm{s}^1`$. The observed annihilation rate will be compared with this generic cross section, which basically only depends on the expansion rate of the universe, i.e. on the value of the Hubble constant. However, it should be noted that this cross section may be energy dependent and the annihilation cross section in the present universe may be much smaller than the value deduced from the time of freeze out, when the temperature was $`m_\chi /22`$ several GeV. On the other hand the annihilation rate may be enhanced by the clustering of DM in “microhaloes”, which increases the density locally. This unknown enhancement factor, usually called ”boost factor”, may vary from a few to a few thousand.$`^{\mathrm{?},\mathrm{?}}`$
## 3 Indirect Dark Matter Detection
The neutral particles play a very special role for indirect DM searches, since they point back to the source. The charged particles change their direction by the interstellar magnetic fields, energy losses and scattering. Therefore the gamma rays provide a perfect means to reconstruct the intensity (halo) profile of the DM by observing the intensity of the gamma ray emissions in the various sky directions. Of course, this assumes that one can distinguish the gamma rays from DM annihilation from the background, mainly from proton-proton interactions. Both for DMA and pp collisions the gamma rays originate mainly from the decay of neutral pions, a light particle produced abundantly in the hadronization process of quarks into hadrons. However, the protons in the galaxies and consequently the quarks inside the protons have a steeply falling energy spectrum ($`NE^{2.7}`$). In contrast, the quarks from DM annihilation are mono-energetic, since the WIMPS annihilate almost at rest, so their mass is converted completely into kinetic energy of the much lighter quarks. Each quark thus obtains an energy corresponding to the mass of the WIMP, which yields a gamma ray spectrum with a sharp cut-off at the mass of the WIMP. So from the shape of the spectrum the WIMP mass can be deduced. The difference in spectral shape between DMA and background allows to obtain their absolute normalizations by fitting their shapes to the EGRET data. These shapes are well known from accelerator experiments and can be obtained e.g. from the PYTHIA code for quark fragmentation $`^\mathrm{?}`$; the parameters in this code have been optimized to fit a wide variety of accelerator data with a single model, the string fragmentation model. The fit of the normalizations can be repeated in many different sky direction to obtain the halo profile of the DM. Given the WIMP number density in all directions from the flux of the excess and the WIMP mass from the spectrum allows to reconstruct the DM mass distribution in our galaxy, which in turn can be used to reconstruct the rotation curve.
A very detailed gamma ray distribution over the whole sky was obtained by the Energetic Gamma Ray Emission Telescope EGRET, one of the four instruments on the Compton Gamma Ray Observatory CGRO, which collected data during nine years, from 1991 to 2000. The EGRET telescope was carefully calibrated in the energy range of 0.1 to 30 GeV, but using Monte Carlo simulations the energy range was recently extended up to 120 GeV $`^\mathrm{?}`$ with a correspondingly larger uncertainty, mainly from the self-vetoing of the detector by the back-scattering from the electromagnetic calorimeter into the veto counters for high energetic showers. It was already noticed in 1997 that the EGRET data showed an excess of gamma ray fluxes for energies above 1 GeV if compared with conventional galactic models.$`^\mathrm{?}`$
Fitting the three contributions of galactic background, extragalactic background and DMA to the energy spectra of 180 independent sky directions yielded astonishingly good fits with the free normalization of the background agreeing reasonably well with the absolute predictions of the galactic models $`^{\mathrm{?},\mathrm{?}}`$ for the energies between 0.1 and 0.5 GeV. Above these energies a clear contribution from Dark Matter annihilation is needed, but the excess in different sky directions can be explained by a single WIMP mass. The fits for 3 different sky directions are shown in Fig. 1.
Alternative explanations for the excess have been plentiful. Among them: weak point sources, which could not be resolved from the background by the EGRET satellite. This is unlikely, since the point sources usually have a rather soft spectrum. If one assumes that most of the unresolved point sources would have similar spectra, their subtraction would reduce the observed diffuse spectra below 1 GeV, but the data above 1 GeV would be much less affected. With our fitting procedure of the shapes, the background is determined by the data below 1 GeV and would thus become lower with unresolved point sources subtracted. Thus would lead to an even stronger excess!
Other ways to increase the excess would be to harden the spectra of the primary nuclei and electrons with respect to the locally measured spectra. Inhomogeneities in the spectra could happen e.g. by density fluctuations from the spiral arms or Supernovae explosions. A summary of these discussions has been given by Strong et al..$`^\mathrm{?}`$ They find that by modifying the electron and proton spectra simultaneously, they can improve the description of the data. However, above 2 GeV the predicted flux of this so-called “optimized” model is still too low, as shown in Fig. 9 of their paper. Since they tried to predict the absolute flux, the overall normalization errors are plotted. However, if one only considers the shape of the spectra, then only the relative systematic errors between the energy points play a role and these are at least a factor two smaller. In this case the probability of the fit, if the shape of the optimized model is fitted to all sky directions, is below $`10^7`$.$`^\mathrm{?}`$ Adding DM to the optimized model improves the fit probability to 0.8 $`^\mathrm{?}`$, of course with a lower boost factor (about factor three), but still a need for DM is evident. Similar results are obtained for the shape proposed by Kamae et al..$`^\mathrm{?}`$ Here the reduction of the boost factor is considerably less, mainly because these authors try to improve the fit by changing the proton spectra only, while in the optimized model both the electron spectra and proton spectra are modified.
An alternative way of formulating the problems of the models without DMA: if the shape of the EGRET excess can be explained perfectly in all sky directions by a gamma contribution originating from the fragmentation of mono-energetic quarks, it is very difficult to replace such a contribution by an excess from nuclei (quarks) (or electrons) with a steeply falling energy spectrum.
From the excess in the various sky directions one can obtain the halo profile under the assumption that the clustering of the DM is similar in all sky directions. This is not necessarily true, since near the centre of the galaxy clumps may be tidally disrupted by the flyby of stars. The annihilation rate is in general proportional to $`B\rho ^n`$, where B is the boost factor and $`n`$ is between 1 and 2, depending on how much of the DM is clustered (n=2 for no clustering and n=1 if all DM is in clusters). Since the EGRET excess measures only the product $`B\rho ^n`$, several choices can be made. For definiteness we use $`n=2`$ and $`B`$ to be the same for all directions and an isothermal halo, which falls like $`1/r^2`$ as expected for a flat rotation curve. The result is surprising: in addition to the isothermal profile the EGRET excess show a substructure in the form of toroidal rings at 4 and 14 kpc, as shown in Fig. 2: on the left hand side the contribution from the $`1/r^2`$ profile is shown, while for the right hand side the ring structure is added. Such enhanced gamma radiation at 4 and 14 kpc was already observed in the original paper on the EGRET excess.$`^\mathrm{?}`$ Note that the appearance of substructure would also be obtained if a radial dependence of $`n`$ and $`B`$ would have been taken. The analysis is sensitive to the radii of ringlike structures, since we are not located at the centre: assuming a constant flux along the ring yields automatically more flux from the nearest parts. The need for these additional rings is most easily seen by comparing the longitudinal profiles in the galactic plane and towards the galactic poles. As shown in Fig. 3 the pole regions are described reasonably well without rings, but for the galactic plane the $`1/r^2`$ profile only describes the data towards the centre. For the larger latitudes one needs the rings, as indicated by the right top panel. Note that for each bin only the flux integrated for data above 0.5 GeV has been plotted.
The position and shape of the outer ring coincides with the ring of stars, discovered in 2003 by several groups.$`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$ These stars show a much smaller velocity dispersion (10-30 km/s) and larger z-distribution than the thick disc, so it cannot be considered an extension of the disc. A viable alternative is the infall of a dwarf galaxy $`^{\mathrm{?},\mathrm{?}}`$, for which one expects in addition to he visible stars a DM component. From the size of the ring and its peak density one can estimate the amount of DM in the outer ring to be $`10^{10}10^{11}`$ solar masses. Since the gamma ray excess requires the full 360 of the sky, one can extrapolate the observed $`100^{}`$ of visible stars to obtain a total mass of $`10^810^9`$ solar masses $`^{\mathrm{?},\mathrm{?}}`$, so the baryonic matter in the outer ring is only a small fraction of its total mass.
The inner ring at 4.2 kpc with a width of 2.1 kpc in radius and 0.2 kpc in $`z`$ is more difficult to interpret, since the density of the inner region is modified by adiabatic compression and interactions between the bar and the halo. However, it is interesting to note that its coordinates coincide with the ring of cold dense molecular hydrogen gas, which reaches a maximum density at 4 kpc and has a width of 2 kpc as well.$`^\mathrm{?}`$ Molecules form from atomic hydrogen in the presence of dust or heavy nuclei. So a ring of neutral hydrogen suggests an attractive gravitational potential in this region, in agreement with the EGRET excess.
To prove that the enhanced gamma ray density is indeed connected to non-baryonic mass the rotation curve was reconstructed from the excess of the diffuse gamma rays in the following way: since the flux determines the number density of DM for a given boost factor and since the mass of each WIMP is between 50 and 100 GeV, one can determine the relative masses of the components (rings plus spherical part) and consequently predict the shape of the rotation curve. The absolute value of the mass can be obtained by requiring that the rotation speed of the solar system is 220 km/s at 8.5 kpc. The two ring model describes the peculiar change of slope at 11 kpc well, as shown in Fig. 4. The contributions from each of the mass terms have been shown separately. The basic explanation for the negative contribution from the outer ring is that a tracer star at the inside of the ring at 14 kpc feels an outward force from the ring, thus a negative contribution to the rotation velocity. It has often been argued that the outer rotation curve cannot be taken seriously, because the errors are large due to the fact that the absolute values of the rotation velocities strongly depend on the value of $`R_0`$, the distance between the solar system and the galactic centre. This is true, as shown by Honma and Sofue$`^\mathrm{?}`$, but they show that the change in slope at about 1.3$`R_0`$ is independent of $`R_0`$. In addition, it has been argued that the inner and outer rotation curve are difficult to compare, since the methods are completely different. The methods are indeed different, but the first 3 data points from the outer rotation curve (between 8 and 11 kpc) show the same slope as the ones from the inner rotation curve, so there seems to be no systematic effect related to the different methods.
## 4 Summary and Outlook
In summary, the EGRET data shows an intriguing hint of DM annihilation, since it explains many unrelated facts simultaneously:
a) An excess of diffuse galactic gamma rays which shows a spectrum consistent with the expectation from WIMP annihilation into gamma rays originating from the fragmentation of mono-energetic quarks.
b) The excess is present in all sky directions with the same spectrum, thus excluding that it originates from anomalous contributions in the centre of the galaxy.
c) The excess shows an strongly increased intensity at positions where extra DM is expected, namely at two doughnut shaped structures at radii of 14 and 4 kpc from the centre of the galaxy. At 14 kpc one has observed a ring of stars thought to originate from the infall of a dwarf galaxy, while at 4 kpc one finds an enhanced concentration of molecular hydrogen thought to form from atomic hydrogen in the presence of dust or heavy nuclei, which can be collected in the gravitational potential of a ring of DM.
d) The enhanced excess of gamma rays cannot be due to additional gas in these rings as proven by the rotation curve calculated from the gamma ray excess: the mass in the rings perfectly describe the hitherto unexplained change of slope in the rotation curve at a distance of about 11 kpc. The amount of visible matter is far too low to have such an impact on the rotation curve.
The results mentioned above make no assumption on the nature of the Dark Matter, except that its annihilation produces hard gamma rays consistent with the fragmentation of monoenergetic quarks between 50 and 100 GeV. WIMPs produce such monoenergetic quarks with energies equal to the WIMP mass. WIMP masses in this range and the observed WIMP self annihilation cross section are consistent with the Lightest Supersymmetric Particle predicted in the Minimal Supersymmetric Model with supergravity inspired symmetry breaking, called the mSUGRA model, if one assumes the enhancement of the annihilation by the clustering of DM to be of the order of 50, which is an order of magnitude not unexpected.$`^{\mathrm{?},\mathrm{?}}`$
Within this supersymmetric model one finds a spin-independent cross section for elastic scattering of a WIMP on a proton of about $`10^{43}\mathrm{cm}^2`$, which is within reach$`^\mathrm{?}`$ of future experiments as shown in Fig. 5. This elastic scattering cross section was calculated with Darksusy$`^\mathrm{?}`$.
Direct and indirect detection experiments do not prove the supersymmetric nature of the WIMPs. If the WIMPs are indeed the lightest supersymmetric particle, then this will become clear at the future LHC collider under construction at CERN in Geneva, where supersymmetric particles of the mass range deduced from the EGRET data$`^\mathrm{?}`$ should be observable from 2008 onwards, if they exist.
In our analysis we only fit the known spectral shapes of the various processes with arbitrary normalizations, so the analysis becomes largely model independent. Interestingly, the normalization factors come out to be in agreement with expectations, both for the WIMP signal and the background.
Alternative models for the EGRET excess without DM have to assume that the locally measured fluxes of protons and electrons are not representative for our galaxy. These models provide significantly worse fits to the data, if one takes the strong correlations in the errors between the different energy bins into account. Of course such models do not explain the stability of the ring of stars at 14 kpc and the change of slope in the rotation curve at $`r=1.3R_0`$.
Therefore the statistical significance of the EGRET excess of at least 10 $`\sigma `$, if fitted to the shape of the diffuse gamma ray background only, combined with all features mentioned above provides an intriguing hint that this excess is indeed indirect evidence for Dark Matter annihilation.
## 5 Acknowledgements
I thank my close collaborators A. Gladyshev, D. Kazakov, C. Sander and V. Zhukov for their contributions to this interesting project. Furthermore I thank V. Moskalenko, A. Strong and O. Reimer for numerous discussions on galactic gamma rays and analysis of EGRET data. This work was supported by the DLR (Deutsches Zentrum für Luft- und Raumfahrt) and a grant from the DFG (Deutsche Forschungsgemeinschaft, Grant 436 RUS 113/626/0-1).
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# Yang-Mills theory in Landau gauge as a liquid crystal
(June 12, 2005)
## Abstract
Using a spin–charge separation of the gluon field in the Landau gauge we show that the $`SU(2)`$ Yang-Mills theory in the low-temperature phase can be considered as a nematic liquid crystal. The ground state of the nematic crystal is characterized by the $`A^2`$ condensate of the gluon field. The liquid crystal possesses various topological defects (instantons, monopoles and vortices) which are suggested to play a role in non-perturbative features of the theory.
preprint: ITEP-LAT/2005-09
Separation of degrees of freedom is a useful analytical tool which is widely used in various physical applications. For example, the spin-charge decomposition (often referred to as the slave-boson formalism ref:slave-boson ) of the strongly correlated electrons is a popular technique invoked to describe a low-temperature physics of the high-$`T_c`$ cuprate superconductors ref:highTc . According to the slave-boson formalism the electron creation operator $`c_{i\sigma }^{}`$ is represented as the product of two operators,
$`c_{i\sigma }^{}=f_{i\sigma }^{}b_i,`$ (1)
where $`i`$ is the lattice site and $`\sigma =,`$ is the spin index. The operator $`f_{i\sigma }^{}`$ creates a chargeless spin-1/2 fermion state (”spinon”) while the operator $`b_i`$ annihilates a charged spin-0 boson state (”holon”). Physically, the electron is represented as a composite of the spinon particle (which carries information about the spin of the electron) and the holon particle (which knows about the electron charge). The decomposition conserves the total number of the degrees of freedom because of the constraint $`f_i^{}f_i+f_i^{}f_i+b_i^{}b_i=1`$. In Eq. (1) the states with double occupancy are disregarded for simplicity.
The local nature of the spin-charge decomposition (1) leads to an emergence of an internal compact $`U(1)`$ gauge symmetry realized in the form of the gauge transformations
$`f_{i\sigma }e^{i\phi _i}f_{i\sigma },b_ie^{i\phi _i}b_i.`$ (2)
Certain properties of the high-$`T_c`$ superconductors can be described ref:highTc ; ref:high\_Tc-U1 by $`U(1)`$ gauge models which are utilizing the mentioned internal gauge symmetry. These gauge models are treatable within the mean field approach which predicts a rich phase structure. In particular, the $`d`$-wave high–$`T_c`$ superconductor is suggested ref:high\_Tc-U1 to be realized as a phase in which the spinon pairing, $`\mathrm{\Delta }_{ij}f_i^{}f_jf_i^{}f_j0`$, is accompanied with a spontaneous breaking of the internal $`U(1)`$ symmetry by the holon condensate:
$`bb_i0.`$ (3)
The presence and the subsequent spontaneous breaking of the internal gauge symmetry may have important physical consequences if even this symmetry is not realized in the original formulation of the theory.
Quantum Chromodynamics is another example of a strongly interacting system in which the breaking of the internal symmetry may play an essential role. Long time ago it was suggested ref:dual:superconductor that the confinement of quarks into hadrons may happen due to a condensation of special gluonic configurations called Abelian monopoles. In this approach – referred to as the dual superconductivity scenario – a condensate of the monopoles breaks spontaneously an internal (or, ”dual”) $`U(1)`$ gauge symmetry. According to the dual superconductor idea, the breaking of the dual symmetry gives rise naturally to the dual Meissner effect, which insures a formation of a QCD string, which in turn leads to the confinement the quarks into hadronic bound states.
The problem of an explicit realization of the dual superconductivity in QCD in terms of the original (gluon) fields is not solved yet. Moreover, the dual superconductivity is shown numerically to be realized ref:dual:superconductivity:numerical only in a special Maximal Abelian gauge which explicitly selects predefined direction(s) in the color gauge group. In this gauge the gluons from the diagonal (Cartan) subgroup are likely to be responsible for the infrared phenomena such as the quark confinement ref:dual:superconductivity:numerical . The off-diagonal gluons were shown to be short–ranged and are largely inessential for the infrared physics ref:off:diagonal .
Recently, it was proposed ref:splitting:niemi to split the gluons in a manner of the spin-charge separation used in the high-$`T_c`$ superconductivity models. The splitting is based on the field decomposition ref:splitting:faddeev which is applied to the off-diagonal gluons while leaving the diagonal gluons intact. In the SU(2) Yang-Mills (YM) theory the splitting of the off-diagonal gluons ref:splitting:niemi ; ref:splitting:faddeev ,
$`A_\mu ^1+iA_\mu ^2=\psi _1𝐞_\mu +\psi _2^{}𝐞_\mu ^{},𝐞_\mu 𝐞_\mu =0,𝐞_\mu 𝐞_\mu ^{}=1,`$ (4)
leads to appearance of two electrically charged (with respect to the Cartan subgroup of the color gauge group) Abelian scalar fields $`\psi _{1,2}`$ and the electrically neutral field $`𝐞_\mu `$ which is a complex vector. There are also other popular gluon field decompositions ref:splitting:cho , some of which were suggested to be related to the monopole condensation.
In this paper we propose a novel generalization of the spin-charge decomposition of the high–$`T_c`$ superconductors (1) to the $`SU(2)`$ Yang–Mills (YM) theory. This decomposition splits the $`SU(2)`$ gluon field into spin and color degrees of freedom treating all color components equally:
$`A_\mu ^a(x)=\mathrm{\Phi }^{ai}(x)e_\mu ^i(x).`$ (5)
Here $`\mathrm{\Phi }^{ai}(x)`$ is the $`3\times 3`$ matrix, and $`a_\mu ^i(x)`$ are the three vectors forming an (incomplete) orthonormal basis in the four dimensional space-time, $`e_\mu ^i(x)e_\mu ^j(x)=\delta ^{ij}`$. The elements of $`\mathrm{\Phi }^{ai}(x)`$ and $`a_\mu ^i(x)`$ are real functions labeled by the color ($`a=1,2,3`$), internal ($`i=1,2,3`$) and Euclidean vector ($`\mu =1,\mathrm{},4`$) indices. Obviously, Eq. (4) is a color-symmetric generalization of Eq. (4). In order to avoid cluttering of notations with lower and upper indices we prefer to work in the Euclidean space–time.
The splitting (5) of the gluon fields can obviously be written in any SU(2) gauge. However, under the local $`SU(2)`$ color transformations, $`A_\mu (x)A_\mu ^\mathrm{\Omega }(x)=\mathrm{\Omega }\left(A_\mu +ig_\mu \right)\mathrm{\Omega }^{}`$, the fields $`\mathrm{\Phi }^{ai}`$ and $`e_\mu ^i`$ mix with each other in a complicated way. Here $`A_\mu A_\mu ^a\sigma ^a/2`$ is the gauge field, $`\sigma ^a`$ are the Pauli matrices, and $`g`$ is the gauge coupling.
In order to make the splitting (5) well–defined we fix the Landau gauge (6) which minimizes the gauge fixing functional,
$`\underset{\mathrm{\Omega }}{\mathrm{min}}F[A^\mathrm{\Omega }],F[A]={\displaystyle \mathrm{d}^4x\left[A_\mu ^a(x)\right]^2},`$ (6)
over the gauge transformations. This gauge condition fixes the $`SU(2)`$ color gauge freedom up to the $`SU(2)`$ global freedom (which is usually disregarded): $`A_\mu ^a(x)\mathrm{\Omega }_{\mathrm{gl}}^{ab}A_\mu ^b(x)`$, where $`\mathrm{\Omega }_{\mathrm{gl}}^{ab}=\mathrm{Tr}\left(\sigma ^a\mathrm{\Omega }\sigma ^b\mathrm{\Omega }^{}\right)/2`$ is the coordinate-independent matrix belonging to the adjoint representation of the color $`SU(2)`$ group.
The transformation rules for the components of the gauge field (5) are:
$`\mathrm{\Phi }(x)\mathrm{\Omega }_{\mathrm{gl}}\mathrm{\Phi }(x)\mathrm{\Lambda }^T(x),e_\mu (x)\mathrm{\Lambda }(x)e_\nu (x)\xi _{\mu \nu },`$ (7)
or, explicitly, $`\mathrm{\Phi }^{ai}(x)\mathrm{\Omega }_{\mathrm{gl}}^{ab}\mathrm{\Phi }^{bj}(x)\mathrm{\Lambda }^{ij}(x)`$ and $`e_\mu ^i(x)\mathrm{\Lambda }^{ij}(x)e_\nu ^j(x)\xi _{\mu \nu }`$. Here $`\mathrm{\Omega }_{\mathrm{gl}}`$ is the matrix of the SU(2) color global transformations, $`\xi _{\mu \nu }`$ is the SO(4) element of the rotations in the Euclidean space-time and $`\mathrm{\Lambda }`$ is the matrix of the internal SO(3) transformations ($`\mathrm{\Lambda }^T\mathrm{\Lambda }=1\mathrm{l}`$):
$`\mathrm{\Omega }_{\mathrm{gl}}SO(3)_{\mathrm{global}}^{\mathrm{color}},\xi SO(4)_{\mathrm{global}}^{\mathrm{spin}},\mathrm{\Lambda }(x)SO(3)_{\mathrm{local}}^{\mathrm{internal}}.`$ (8)
Equation (5) can be interpreted as a spin-color separation of the gluon field since the color gauge transformations $`\mathrm{\Omega }_{\mathrm{gl}}`$ are acting only on the matrix field $`\mathrm{\Phi }`$ while spin transformations $`\xi `$ affect only the field $`e_\mu `$. Note that the color and the space-time rotations are the symmetries of the original $`SU(2)`$ gauge theory while the internal symmetry group originates from the form of the splitting (5) in analogy to the compact $`U(1)`$ internal symmetry (2) of the high-$`T_c`$ superconductors. For the sake of brevity we call below the group of the internal gauge transformations as $`SO(3)_{\mathrm{int}}`$.
The proposed splitting (5) is self-consistent from the point of view of counting of the degrees of freedom (d.o.f.). The original field $`A_\mu ^a`$ is described by $`3\times 4=12`$ real functions<sup>1</sup><sup>1</sup>1While counting the degrees of freedom we do not take into account the pure color gauge degrees of freedom ($`3`$ d.o.f.) and do not impose the Gauss constraint ($`3`$ d.o.f.) to select the physical states because these restrictions equally affect both sides of Eq. (5).. The field $`A_\mu ^a`$ is now rewritten (5) as the product of the matrix $`\mathrm{\Phi }^{ai}`$ ($`3\times 3=9`$ d.o.f.) and the three vector fields $`e_\mu ^i`$ ($`3\times 4=12`$ d.o.f.) subjected to the orthonormality constraints ($`6`$ d.o.f.). The group $`SO(3)_{\mathrm{int}}`$ of the internal gauge transformations has 3 generators ($`3`$ d.o.f.). Thus, the number of the d.o.f. in the field $`A_\mu ^a`$ (which is 12) is the same as the total number of d.o.f. in the product of the fields $`\mathrm{\Phi }^{ai}`$ and $`e_\mu ^i`$: (which is $`9+1263=12`$).
It is instructive to rewrite the SU(2) gauge model in an explicitly $`SO(3)_{\mathrm{int}}`$ invariant form. To this end one may introduce two composite gauge fields:
$`\mathrm{\Gamma }_\mu ^{ij}={\displaystyle \frac{1}{2}}\left(e_\nu ^i_\mu e_\nu ^je_\nu ^j_\mu e_\nu ^i\right),\vartheta _\mu ^{ij}={\displaystyle \frac{1}{2}}\left((\mathrm{\Phi }^1)^{ia}_\mu \mathrm{\Phi }^{aj}(\mathrm{\Phi }^1)^{ja}_\mu \mathrm{\Phi }^{ai}\right),`$ (9)
and two composite matter fields,
$`\chi ^{ij}=\mathrm{\Phi }^{ai}\mathrm{\Phi }^{aj},z_\mu ^{ij}={\displaystyle \frac{1}{2}}\left((\mathrm{\Phi }^1)^{ia}_\mu \mathrm{\Phi }^{aj}+(\mathrm{\Phi }^1)^{ja}_\mu \mathrm{\Phi }^{ai}\right),`$ (10)
which transform under the internal gauge transformations as
$`\mathrm{\Gamma }_\mu \mathrm{\Lambda }(\mathrm{\Gamma }_\mu +_\mu )\mathrm{\Lambda }^T,\vartheta _\mu \mathrm{\Lambda }(\vartheta _\mu +_\mu )\mathrm{\Lambda }^T,\chi \mathrm{\Lambda }\chi \mathrm{\Lambda }^T,z_\mu \mathrm{\Lambda }z_\mu \mathrm{\Lambda }^T.`$ (11)
The $`SO(3)_{\mathrm{int}}`$ gauge fields $`\mathrm{\Gamma }_\mu `$ and $`\vartheta _\mu `$ are asymmetric with respect to permutations of the internal indices while the scalar matter field $`\chi `$ and the vector matter field $`z_\mu `$ are symmetric under these permutations. The matter fields transform in the adjoint representation of the $`SO(3)_{\mathrm{int}}`$ gauge group. Note that it is impossible to construct composite matter fields from the ”spin” field $`e_\mu `$ in a manner of Eq. (10) due to the orthonormality constraints imposed on $`e_\mu `$.
The Landau gauge functional (6) can be expressed in terms of the matter field $`\chi `$
$`F[A][\chi ]={\displaystyle \mathrm{d}^4x\mathrm{Tr}\chi }.`$ (12)
Note that this functional still invariant under all global and local transformations (8).
Technically, the existence of the two gauge fields (9) and one adjoint vector field (10) allows us to define an arbitrary number of covariant derivatives, $`D_\mu ^{ij}(\gamma )=_\mu \delta ^{ij}+\gamma _\mu ^{ij}`$ where the vector field $`\gamma _\mu `$ stands for any linear combination of the $`\mathrm{\Gamma }_\mu `$, $`\vartheta _\mu `$ and $`z_\mu `$ fields which transforms as a $`SO(3)_{\mathrm{int}}`$ gauge field. Then, the derivative of the gauge field $`A_\mu ^a`$ can be represented in an explicitly $`SO(3)_{\mathrm{int}}`$ invariant form, $`_\mu A_\nu ^a=(\widehat{\mathrm{\Phi }}^a,D_\mu \widehat{e}_\nu )+(D_\mu \widehat{\mathrm{\Phi }}^a,\widehat{e}_\nu )`$. The local (differential) condition of the Landau gauge, $`_\mu A_\mu ^a=0`$, can be rewritten as a constraint
$`(\widehat{\mathrm{\Phi }}^a,D_\mu \widehat{e}_\mu )+(D_\mu \widehat{\mathrm{\Phi }}^a,\widehat{e}_\mu )=0.`$ (13)
Here vectors $`\widehat{\mathrm{\Phi }}^a(\mathrm{\Phi }^{a1},\mathrm{\Phi }^{a2},\mathrm{\Phi }^{a3})^T`$ are the columns of the matrix $`\mathrm{\Phi }^{ai}`$, $`\widehat{e}_\mu =(e_\mu ^1,e_\mu ^2,e_\mu ^3)^T`$, and $`(a,b)=a^ib^i`$ is the scalar product in the internal $`SO(3)_{\mathrm{int}}`$ space. Below we make the choice $`\gamma _\mu ^{ij}=\mathrm{\Gamma }_\mu ^{ij}`$ for convenience.
It is also convenient to introduce the vector $`e_\mu ^4=\epsilon _{\mu \nu \alpha \beta }e_\nu ^1e_\alpha ^2e_\beta ^3`$. The four vectors $`e_\mu ^{\overline{i}}`$, $`\overline{i}=1,\mathrm{},4`$ form a complete orthonormal basis in the $`4D`$ space-time, $`e_\mu ^{\overline{i}}e_\mu ^{\overline{j}}=\delta ^{\overline{i}\overline{j}}`$. The internal $`SO(3)_{\mathrm{int}}`$ transformations act in the subspace spanned onto vectors $`e_\mu ^k`$ with $`k=1,2,3`$ while leaving the vector $`e_\mu ^4`$ intact.
The YM Lagrangian be divided into the three parts
$`L_{SU(2)}[A]{\displaystyle \frac{1}{4}}\left[G_{\mu \nu }^a(A)\right]^2=L_0[\mathrm{\Phi },\chi ,\mathrm{\Gamma }]+L_1[\chi ,\mathrm{\Gamma },\vartheta ]+L_2[\chi ]+L_{\mathrm{gf}},`$ (14)
where $`G_{\mu \nu }^a(A)=_{[\mu ,}A_{\nu ]}^a+g\epsilon ^{abc}A_\mu ^bA_\nu ^c`$ is the $`SU(2)`$ field strength tensor and the term $`L_n`$ is proportional to the $`n^{\mathrm{th}}`$ power of the $`SU(2)`$ coupling constant $`g`$. For a moment we disregard the term $`L_{\mathrm{gf}}`$ coming from the Landau gauge fixing. Using an appropriate multiplication by the vectors $`e_\mu ^{\overline{k}}`$ to convert the Euclidean indices into the internal $`SO(3)_{\mathrm{int}}`$ basis we rewrite the YM Lagrangian (14) as follows:
$`L_0[\mathrm{\Phi },\chi ,\mathrm{\Gamma }]`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(𝒟_{\overline{k}}(\mathrm{\Gamma })\widehat{\mathrm{\Phi }}^a\right)^2+{\displaystyle \frac{1}{2}}(\mathrm{\Sigma }(\mathrm{\Gamma }),\chi \mathrm{\Sigma }(\mathrm{\Gamma })),`$ (15)
$`L_1[\chi ,\mathrm{\Gamma },\vartheta ]`$ $`=`$ $`2g\sqrt{\mathrm{det}\chi }(\mathrm{\Gamma }_k^{ij}\vartheta _k^{ij})\epsilon _{ijk},`$ (16)
$`L_2[\chi ]`$ $`=`$ $`{\displaystyle \frac{g^2}{4}}\left[\left(\mathrm{Tr}\chi \right)^2\mathrm{Tr}\chi ^2\right].`$ (17)
where $`𝒟_{\overline{k}}(\mathrm{\Gamma })e_\mu ^{\overline{k}}D_\mu (\mathrm{\Gamma })`$ is the covariant derivative acting on the internal $`SO(3)_{\mathrm{int}}`$ indices. Note that the spin field $`\widehat{e}_\mu `$ enters the Lagrangian (14) only in the form of the connection $`\mathrm{\Gamma }_{\overline{k}}^{ij}`$.
In order to simplify the $`L_0`$ part of the YM Lagrangian (15) we used the differential Landau gauge condition and neglected a full-derivative surface term. The first term in $`L_0`$ is the kinetic term for the ”color” component of the gluon field $`\widehat{\mathrm{\Phi }}^a`$ in the background of the $`SO(3)_{\mathrm{int}}`$ gauge field $`\mathrm{\Gamma }`$. The second term in $`L_0`$ can be interpreted as a ”dielectric” energy density associated with the (space-dependent) ”dielectric susceptibility” $`\chi `$ and the (dynamical) $`SO(3)_{\mathrm{int}}`$ ”electric field” $`\mathrm{\Sigma }^i(x)=\mathrm{\Lambda }_{}^{ij}(x)_j(x)`$. Here the $`SO(3)_{\mathrm{int}}`$ gauge transformation $`\mathrm{\Lambda }_{}^{ij}`$ diagonalizes the matrix $`\mathrm{\Gamma }_{\overline{k}}^{4i}\mathrm{\Gamma }_{\overline{k}}^{4j}=[\mathrm{\Lambda }_{}\mathrm{diag}(_1^2,_2^2,_3^2)\mathrm{\Lambda }_{}^T]^{ij}`$ with $`_i(x)0`$.
The second part (16) of the Lagrangian represents the interaction between the gauge fields $`\mathrm{\Gamma }`$ and $`\vartheta `$ with the effective coupling $`g\mathrm{det}^{1/2}\chi g\mathrm{det}\mathrm{\Phi }`$. The third part (17) is a local potential $`V(\chi )`$ on the ”dielectric susceptibility” field $`\chi `$.
The analogy of the spin-color separation of the gluon in YM theory (5) with the spin-charge separation of the electron in the high-$`T_c`$ superconductor models ref:highTc manifests itself also in the absence of the kinetic terms for the composite gauge fields $`\mathrm{\Gamma }_\mu `$ and $`\vartheta _\mu `$. This fact is natural since the local construction of each of the composite gauge fields (9) involves already a single derivative while canonical local Lagrangians (i.e., the YM Lagrangian) contain terms with at most two derivatives. The only explicitly propagating field in formulation (14) is $`\widehat{\mathrm{\Phi }}^a`$.
Besides the remarkable analogy of the spin–color separation in the YM theory with the spin-charge separation in the high–$`T_c`$ superconductivity, the YM theory has another interesting analogue in the condensed matter physics. Namely, the YM Lagrangian (14-17) can be interpreted as the free energy density of a nematic liquid crystal.
The ordinary nematic crystals ref:nematics:review consist of rod-like molecules which tend to align parallel to a direction $`𝐧(𝐱,t)`$. The molecule is invariant under reflections with respect to a plane perpendicular to the molecule axis. The unit vector $`𝐧`$ – called the Frank director – is chosen spontaneously in the absence of external electric or magnetic fields. The molecules in liquid crystals do not have a positional order contrary to solid crystals characterized by lattice-like structures. The energetically favored ground state of the nematic crystal is realized at low temperatures and is characterized by a constant director field, $`𝐧(𝐱,t)=𝐧_0`$. As temperature increases the system undergoes a transition from the nematic phase to the ordinary (isotropic) phase.
Due to the symmetries of the nematic molecule the symmetry group of the ordinary nematic is $`G=SO(3)/ZZ_2`$. Therefore, the order parameter in a nematic may be a unit vector but without associated direction ref:nematics:review (i.e., a vector without arrowhead). However, it is more convenient to define the order parameter to be diadic in $`n_i`$ similarly to the diamagnetic (or, dielectric) susceptibility $`\stackrel{~}{\chi }_{\alpha \beta }`$. The excellent candidate for the order parameter which discriminates between the nematic and isotropic phases ref:nematics:review is the amount of disorder in $`\stackrel{~}{\chi }_{\alpha \beta }`$:
$`\stackrel{~}{Q}_{\alpha \beta }=\stackrel{~}{\chi }_{\alpha \beta }{\displaystyle \frac{1}{3}}\delta _{\alpha \beta }\stackrel{~}{\chi }_{\gamma \gamma }=\mathrm{\Delta }\stackrel{~}{\chi }{\displaystyle \underset{s}{}}\left(n_\alpha ^{(s)}n_\beta ^{(s)}{\displaystyle \frac{1}{3}}\delta _{\alpha \beta }\right),`$ (18)
where the last equality is written for the molecules with exact axial symmetry. In Eq. (18) the summation is going over all molecules in a small but macroscopic volume, $`𝐧^{(s)}`$ is the direction of the axis of the $`s^{\mathrm{th}}`$ molecule, and $`\mathrm{\Delta }\stackrel{~}{\chi }=\stackrel{~}{\chi }_{}\stackrel{~}{\chi }_{}`$ is the anisotropy in the diamagnetic (dielectric) susceptibility along and perpendicular to the molecule axis. The quantity $`\stackrel{~}{Q}_{\alpha \beta }`$ is non-zero in the nematic phase while it vanishes in the isotropic phase. Below we refer to $`\stackrel{~}{\chi }`$ as to the dielectric susceptibility.
The dependence of the free energy on the order parameter (18) is usually given by an effective Landau–Lifshitz (LL) potential ref:nematics:review ,
$`F_{LL}(\stackrel{~}{Q})=F_0+{\displaystyle \mathrm{d}^3𝐱\underset{n2}{}\alpha _n\mathrm{Tr}\stackrel{~}{Q}^n}`$ (19)
where $`\alpha _n`$ are functions of temperature $`T`$. The dependence of the free energy on the isotropic factor $`\mathrm{Tr}\stackrel{~}{\chi }`$ may be included into the free energy of the normal state, $`F_0`$.
The deviations of the Frank director $`𝐧`$ from the ground state $`𝐧_0`$ are typically described by the Oseen–Zöcher–Frank (OZF) free energy,
$`F_{OZF}[𝐧]={\displaystyle \frac{1}{2}}{\displaystyle \mathrm{d}^3𝐱\left[K_1(𝐧)^2+K_2(𝐧\times 𝐧)^2+K_3(𝐧\times \times 𝐧)^2\right]},`$ (20)
where the first three terms describe the free energy associated with the splay, twist and bend distortions. The total free energy of the nematic crystal is $`F[\stackrel{~}{Q},𝐧]=F_{LL}(\stackrel{~}{Q})+F_{OZF}[𝐧]`$. Note that relation (18) makes it possible to rewrite the OZF free energy as a more complicated (compared to (20)) expression in terms of the order parameter $`\stackrel{~}{Q}`$.
The YM theory (14-17) can be associated with a nematic crystal in which the ”molecules” are directed in the internal $`SO(3)_{\mathrm{int}}`$ space. There are three species of equivalent molecules in each space-time point (the number of species equals to the number of the gluons, $`N_c=3`$). Consequently, the direction of the local color field in the YM theory, $`(\widehat{\mathrm{\Phi }}^a)^i(x)/|\widehat{\mathrm{\Phi }}^a(x)|`$, is associated with the direction $`n_i^{(a)}(x)`$ of the $`a^{\mathrm{th}}`$ molecule species in the point $`x`$. Then the adjoint matter field $`\chi ^{ij}=_a\mathrm{\Phi }^{ai}\mathrm{\Phi }^{aj}`$ can be associated with the dielectric susceptibility, $`\stackrel{~}{\chi }_{\alpha \beta }=\mathrm{\Delta }\stackrel{~}{\chi }_sn_\alpha ^{(s)}n_\beta ^{(s)}`$. Note that YM ”dielectric susceptibility” $`\chi `$ is diadic in the fields $`\widehat{\mathrm{\Phi }}^a`$ similarly to the dielectric susceptibility $`\stackrel{~}{\chi }`$ of the nematic.
The proposed association is largely based on the form of the YM term $`L_2(\chi )`$, Eq. (17), which plays a role of the LL potential (19) for the YM ”dielectric” field $`\chi `$. This term can be rewritten via the isotropic factor $`\mathrm{Tr}\chi `$ and the traceless symmetric matrix $`Q^{ij}`$, constructed from the ”susceptibility” $`\chi ^{ij}`$ similarly to the nematic case (18): $`L_2[\chi ]=\frac{g^2}{6}(\mathrm{Tr}\chi )^2\frac{g^2}{4}\mathrm{Tr}Q^2`$. The negative sign in front of the second term leads to the instability to develop a disorder in the ”dielectric susceptibility” $`\chi ^{ij}`$.
The $`L_0`$ term, Eq. (15), is a covariant generalization of the kinetic part of the OZF free energy (20) corresponding to the liquid crystal whose splay, twist and bend distortion constants are equal, $`K_1=K_2=K_3=1`$. Indeed, in this case the first three terms in Eq. (20) are reduced to $`\frac{1}{2}_{i,j=1}^3(_in_j)`$. Then, we get the $`L_0`$ term in the YM Lagrangian by (i) imposing the natural requirement of the $`SO(3)_{\mathrm{int}}`$ covariance, $`_\mu D_{\overline{k}}(\mathrm{\Gamma })`$, and (ii) taking into account all molecule species, $`𝐧\widehat{\mathrm{\Phi }}^a(x)`$.
As for the $`L_1`$ term, Eq. (16), it can be interpreted as an energy density associated with a mutual non–alignment of the directions of the different molecule species $`(\widehat{\mathrm{\Phi }}^a)^i(x)/|\widehat{\mathrm{\Phi }}^a(x)|`$.
Let us find the ground state of the nematic associated with the YM theory (15,16,17). In terms of the eigenvalues of the matrix $`\chi =\mathrm{diag}(\chi _1,\chi _2,\chi _3)`$, the ground state $`\chi =\chi ^{(0)}`$ is defined by the relations:
$`{\displaystyle \underset{\stackrel{i,j=1}{i>j}}{\overset{3}{}}}\chi _i^{(0)}\chi _j^{(0)}=0,{\displaystyle \underset{i=1}{\overset{3}{}}}\chi _i^{(0)}0,{\displaystyle \underset{i=1}{\overset{3}{}}}\chi _i^{(0)}0,`$ (21)
where the first relation comes from the condition $`\mathrm{Tr}\chi ^2=(\mathrm{Tr}\chi )^2`$ corresponding the global minimum of the Ginzburg–Landau potential (17). The last two relations in Eq. (21) come from the specific definition of the $`\chi `$–field (10) implying that $`\mathrm{Tr}\chi _{ai}(\mathrm{\Phi }^{ai})^20`$ and $`\mathrm{det}\chi (\mathrm{det}\mathrm{\Phi })^20`$, respectively. Equations (21) imply that at least two eigenvalues of $`\chi `$ must be zero. Without loss of generality we take $`\chi _1^{(0)}=\chi _2^{(0)}=0`$, and therefore the ground state is $`\chi ^{(0)}=\mathrm{diag}(0,0,\chi _0)`$, where $`\chi _00`$ is not fixed.
The perturbative vacuum (in terms of the original gluon fields $`A_\mu ^a`$) corresponds to $`\chi _0=0`$, i.e. to the isotropic liquid state. What makes the YM field similar to the nematic liquid is the non–perturbative part of $`\chi _0`$, which is fixed by the minimum of the Landau gauge functional (12). This minimum is nothing but the $`A^2`$–condensate ref:A2:condensate:general , $`\mathrm{Tr}\chi =A_\mu ^2`$, evaluated in the Landau gauge. Thus, the isotropic liquid state is broken to the nematic crystal state by the $`A^2`$ condensate. This spontaneous symmetry breaking of the isotropic $`SO(3)_{\mathrm{int}}`$ is similar to the breaking of the compact gauge group by the holon condensate (3).
Technically, a particular non-zero value of the $`A^2`$-condensate emerges due to the presence of the gauge–fixing term $`L_{\mathrm{gf}}`$ in Eq. (14) which also contributes to the free energy of the nematic liquid and which was disregarded till now. According to the numerical calculations of the $`A^2`$ condensate ref:A2:numerical , $`g^2\chi _0(3\text{GeV})^2`$.
The non-perturbative vacuum state, $`\chi ^{(0)}=\mathrm{diag}(0,0,\chi _0)`$ with $`\chi _0>0`$, is still invariant under the (unbroken) group of rotations about the third axis in the internal space, $`H=SO(2)_{\mathrm{int}}`$. Due to the fact that the $`SO(3)_{\mathrm{int}}`$ gauge field $`\mathrm{\Gamma }`$ is non-propagating, the partial spontaneous breaking of the original internal symmetry does not lead to a massless vector field.
The interesting question is a possible existence of topological defects which are generally characterized by non-trivial homotopic groups $`\pi _n(G/H)`$ of the vacuum manifold $`G/H`$ of the model. The vacuum manifold of the YM theory with Lagrangian written in the form (14-17) is similar to the vacuum manifold of an ordinary nematic ref:nematics:topological with $`G/H=SO(3)/(ZZ_2\times SO(2))`$. In particular, the nematic state contains the $`Z_2`$ vortices since $`\pi _1(G/H)=ZZ_2`$. This feature may make the physics of the YM nematic state similar to the center vortex picture of the quark confinement in the YM theory ref:center:vortex .
Moreover, the nematic crystal contains monopole-like defects characterized by non-negative integers since $`\pi _2(G/H)=ZZ/ZZ_2Z_+=0,1,2,\mathrm{}`$. The monopoles have the hedgehog–like structure constructed from the arrowless ”molecules” (the last fact leads to an identification of the monopoles with anti-monopoles). The presence of the monopoles may provide a relation between the nematic liquid crystal and the dual superconductor in the YM theory ref:dual:superconductor . A signature of this relation may already be found in Ref. ref:suzuki by observing the dual Meissner effect in the Landau gauge. Finally, the third homotopy group of the vacuum manifold is also nontrivial, $`\pi _3(G/H)=ZZ`$, which may have a link to the instanton physics.
The disorder, caused by the presence of the described topological defects in the Landau gauge may lead to the non–trivial consequences for the non–perturbative physics of the YM theory similarly to the effects caused by the center vortex percolation ref:center:vortex and by the Abelian monopole condensation ref:dual:superconductor .
Finally, we note that lattice simulations ref:A2:condensate:temperature indicate that the $`A^2`$ condensate drops by amount of 92% at the finite-temperature phase transition, $`T=T_c`$. Therefore one may expect that in the deconfinement phase, $`T>T_c`$, the $`4D`$ nematic state may transform to a $`3D`$ nematic state characterized by much lower value of the ”nematic dielectric susceptibility” $`\chi `$. Since the spatial dynamics of the gluon fields remains non–perturbative in the deconfinement phase, one may expect that the nematic crystal splits into two modes: the temporal components of the gluon fields form an ordinary ”isotropic liquid” while the spatial components are still in a nematic state.
Summarizing, the spin-charge separation idea – originally invented to describe properties of the high-$`T_c`$ superconductors – may also be applied to the YM theory in the form of the spin-color separation. This approach allows to identify the ground state of the low-temperature phase YM theory in the Landau gauge with a nematic liquid crystal. The perturbative isotropic liquid state is broken down to the nematic liquid crystal state by the $`A^2`$ condensate. The nematic crystal contains various topological defects which may play a role in explaining of non-perturbative features of the YM theory.
###### Acknowledgements.
The author is supported by grants RFBR 04-02-16079 and MK-4019.2004.2. The author is grateful to F.V. Gubarev, A. Niemi and M.I.Polikarpov for useful discussions.
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# X-ray Emission from the Weak-lined T Tauri Binary System KH 15D
## 1 Introduction
KH 15D is a unique eclipsing pre-main sequence (PMS) system near the Cone nebula in NGC 2264 (Kearns et al., 1997; Kearns and Herbst, 1998). The visible star is of K6 or K7 spectral class (Hamilton et al., 2001; Agol et al., 2004) and has an H$`\alpha `$ equivalent width of $``$ 2 Å, typical of a weak-lined T Tauri star (WTTS). Its mass and age are $``$0.6 $`M_{}`$ and 2 My, respectively (Hamilton et al., 2001). At high spectral resolution the star reveals broad wings on its hydrogen emission lines and during eclipse one clearly sees forbidden emission lines (Hamilton et al., 2003). These features signify that accretion and outflow are still active in the system, but not at the level of a typical classical T Tauri star (CTTS). While the star has no measured infrared excess, an apparent disk and jet in H<sub>2</sub> have been detected by Deming, Charbonneau & Harrington (2004) and Tokunaga et al. (2004), respectively.
The extremely long duration of the eclipse, currently about one-half of the period, clearly shows that the eclipsing body is not a companion star. Rather, it appears to be part of a circumstellar or circumbinary disk (Herbst et al., 2002). During eclipse the system becomes both bluer and more highly polarized (Herbst et al., 2002; Agol et al., 2004), suggesting that we are seeing it primarily or entirely in scattered light. There are two time scales associated with the eclipse, a 48.37 day cycle for the main eclipse and a secular increase in the eclipse duration of about 1 day per year (Herbst et al., 2002; Winn et al., 2003; Hamilton, 2004). Two recent models based on the historic light curve (Winn et al., 2003; Johnson & Winn, 2004) have proposed that the 48 day eclipse cycle is the orbital period of a binary system while the secular variation is caused by precession of the circumbinary disk (Winn et al., 2004; Chiang & Murray-Clay, 2004). If this is correct it means that, for the first time, we can probe the structure of a disk on length scales as small or smaller than a stellar diameter and monitor events in a possibly planet-forming disk on human time scales! Clearly it is important to understand as much as possible about this unique PMS stellar system and to exploit its fortuitous geometry while the opportunity lasts.
One characteristic of T Tauri stars, especially WTTS, is that they are prodigious sources of X-ray emission, although for still largely unknown reasons (Feigelson & Montmerle, 1999; Feigelson et al., 2003). We hoped to use the periodic eclipse of the K6-7 star behind an optically thick and presumably X-ray opaque circumstellar disk, to allow us to map the structure of the coronal plasma in this WTTS. As a prelude to this intended study we searched for archival X-ray data on NGC 2264 and found a long exposure in the archives of the Chandra X-ray Observatory that includes KH 15D. It was obtained during a time interval when the star was out of eclipse so we expected a relatively strong signal, characteristic of a WTTS. Instead, we found that the total X-ray count out of eclipse is so small that it may not be possible to learn much by monitoring during an eclipse cycle. This in turn has prompted us to consider the extent to which KH 15D is unusual in yet another way, namely as a remarkably faint X-ray source for a WTTS. In this paper we present the case that it is, indeed, an unusually weak source of X-ray emission and discuss possible implications of this for the system and for the broader question of X-ray production in solar-like, PMS stars.
## 2 X-ray Data
### 2.1 The X-ray Luminosity of KH 15D
The southern portion of NGC 2264, containing KH 15D, was observed with the ACIS-I array on board the Chandra X-ray Observatory for 95.74 ks on 2002 October 28–29 (UT). The instrument was operated in its nominal mode, with a frame time of 3.2 s and a focal-plane temperature of –120 C. We reprocessed the data using the CIAO software, version 3.2.1, to correct for charge-transfer inefficiency of the front-illuminated CCDs and to apply a time-dependent gain correction to the data. In addition, we improved the original afterglow correction and removed the 0.5-pixel randomization applied during the standard processing of the data. The screened data used for analysis consist of events with grades 0, 2, 3, 4, and 6 and energies between 0.5 and 8.0 keV.
The J2000 coordinates of the instrument aimpoint for the observation are $`\alpha `$ = $`06^\mathrm{h}40^\mathrm{m}58.^\mathrm{s}`$10, $`\delta `$ = $`+09^{}34^{}00.{}_{}{}^{\prime \prime }40`$; thus, KH 15D, at $`\alpha `$ = $`06^\mathrm{h}41^\mathrm{m}10.^\mathrm{s}`$27, $`\delta `$ = $`+09^{}28^{}33.{}_{}{}^{\prime \prime }40`$, is included within the $`17^{}\times 17^{}`$ ACIS-I field of view at an off-axis angle of 6$`\stackrel{}{\mathrm{.}}`$35. A weak source at the location of KH 15D is faintly visible in the image. We used bright nearby sources to determine the appropriate size of the source aperture (12 pixels, or $`6^{\prime \prime }`$, in radius), which was centered on the optical position of KH 15D. The background level was estimated in a concentric, source-free annulus with inner and outer radii of 20 and 60 pixels, respectively. A total of 22.5 net counts were detected in the full 0.5–8.0 keV band, corresponding to a signal-to-noise ratio of 3.5. The spectrum of KH 15D is soft, with $`80`$% of the net counts falling below 2 keV. The implied band ratio (i.e., the ratio of counts detected in the hard 2–8 keV and soft 0.5–2 keV ranges) is 0.264, although, because KH 15D is not formally detected in the hard band, this should be considered an upper limit.
We used the XSPEC software, along with the response matrix and effective area files generated by CIAO, to estimate the X-ray flux of KH 15D. A spectral model consisting of a single-temperature, optically thin (MEKAL) plasma was adopted (Feigelson et al., 2002). Solar abundances were assumed for the plasma, and based on the maximum possible reddening of KH 15D of $`E(BV)=0.1`$ (Hamilton et al., 2001), an absorption column density of $`2\times 10^{20}`$ atoms cm<sup>-2</sup> was included in the model. We adjusted the temperature of the plasma until the above band ratio was obtained, which suggests $`kT=2.7`$ keV. Using the model normalization required to match the observed count rate of the source, we obtain a 0.5–8 keV flux of $`2.2\times 10^{15}`$ erg cm<sup>-2</sup> s<sup>-1</sup>. For a distance of 760 pc (Sung, Bessel and Lee, 1997; Hamilton et al., 2001), this corresponds to an unabsorbed X-ray luminosity of $`1.5\times 10^{29}`$ erg s<sup>-1</sup>. Given that the maximum values of the band ratio and column density were assumed, this should be taken as an upper limit.
### 2.2 Comparison with Other PMS Stars
Our derived total X-ray luminosity for KH 15D may be compared with other PMS stars of similar effective temperature in NGC 2264 and in the Orion Nebula Cluster (ONC). The most direct comparison is with stars on the same archival Chandra image as KH 15D. We used the photometric and spectroscopic surveys of Rebull et al. (2002) and Lamm et al. (2004) to search for K stars within that field that had comparable magnitude (within 0.7 mag in $`I`$) and color (within 0.04 mag in $`VI`$). Since the reddening in NGC 2264 is relatively small (E(B-V) $`0.1`$ or less) and uniform, this procedure should result in a reasonable comparison set. Five stars meeting the photometric conditions were found and all five were firmly detected in the Chandra image. Their source counts and X-ray spectra were extracted in the same manner as described above for KH 15D.
Because of the strength of the detections of the comparison stars, we were able to estimate their X-ray fluxes directly via spectral modeling. We began by fitting single-temperature plasma models to each of the spectra. However, good fits were obtained in just two of the cases. For the remaining three objects, we employed two-temperature plasma models, which are frequently required to fit the X-ray spectra of PMS stars (Feigelson et al., 2002). The fitted column densities for four of the five objects are consistent with zero. The fit of other object, star 3748, suggests a column density of $`4\times 10^{21}`$ cm<sup>-2</sup>. This absorption was not corrected for when calculating the star’s X-ray flux, so the X-ray luminosity we derive for it is a lower limit.
Table 1 summarizes the optical and X-ray properties of the five comparison stars and KH 15D. Listed for each object are the source number from Rebull et al. (2002), the $`I`$-band magnitude, the $`VI`$ color, the spectral type, the net counts detected in the 0.5–8 keV band, the hard-to-soft band counts ratio $`H/S`$, the plasma temperature(s) in keV, the 0.5–8 keV flux (in erg cm<sup>-2</sup> s<sup>-1</sup>), and the X-ray luminosity (in erg s<sup>-1</sup>) in the same band. As the table indicates, KH 15D is significantly underluminous relative to this set of comparison stars. Moreover, it is evident that the weakness of its X-ray emission is not a result of a greater amount of soft X-ray absorption: the upper limit to its band ratio and the inferred plasma temperature are consistent with those of the comparison stars. We also calculated X-ray fluxes for the ONC stars using exactly the same spectral model as for KH15D, with resultant fluxes that differed by only 3 to 25% (with an average deviation of 17%) from what is shown in Table 1.
To expand the comparison set, we have also employed observations of K3-K7 stars in the northern part of NGC 2264 whose X-ray luminosities were calculated by Ramirez et al. (2004) based on an ACIS-I Chandra image. They detected 37 likely cluster members in this spectral range, as well as 5 non-members (Rebull & Stauffer, 2005), and their exposure time of 48.1 ks is long enough to have reached sources close to the luminosity of KH 15D, if not below it. The X-ray luminosity of their cluster members is derived in a manner similar to what we have used and is based on the same assumed distance. In Fig. 1 we compare the X-ray luminosity of KH 15D (large square) to the the stars in the northern part of the cluster (crosses) and to those in the southern part from Table 1 (diamonds). It is clear from this comparison that KH 15D has a lower than typical X-ray luminosity for mid-K stars in both parts of NGC 2264. It lies about an order of magnitude below the median value of both samples and is more than a factor of three fainter than the next faintest detected object.
As a further test of the degree to which KH 15D is anomalously weak in X-rays, we compare it to stars of similar spectral class in the ONC. This cluster is slightly younger and about half the distance of NGC 2264. The optical data come from the extensive photometric and spectroscopic survey by Hillenbrand (1997). Total X-ray luminosities in the Chandra band of 0.5–8 keV have been derived by Feigelson et al. (2005) for more than a thousand sources in the ONC, based on an extraordinarily long exposure of 850 ks. Fig. 2 shows a comparison of our result for KH 15D with the 74 ONC members having spectral types between K5 and K8, inclusive. The inferred plasma temperatures of these stars are comparable to the value of 2.7 keV that we infer for KH 15D, which is typical of PMS stars in general (Feigelson et al., 2002). Note that this X-ray survey detected every optically known ONC member in this spectral range, so the comparison sample is as complete as possible. It is clear from this figure that KH 15D is a very weak X-ray source compared to the ONC stars of similar effective temperature.
We further note that KH 15D is likely to be even more anomalously weak as an X-ray emitter than is evident from Fig. 2. There are two reasons for this. First, the extinction in the ONC is much higher in general than in NGC 2264 and is also highly variable. For example, two of the three ONC stars in Fig. 2 with cited X-ray luminosity less than KH 15D have visual extinction estimates exceeding 5 magnitudes! This suggests to us that some of the apparent scatter in X-ray luminosity in the cluster, especially the scatter to low values, is due to extinction, not lack of X-ray production. Second, we note that the ONC is slightly younger and probably has, therefore, a higher percentage of CTTS compared to WTTS. This is hard to verify because of the strong nebulosity in which the ONC is embedded. If true, however, it means that there may be more low luminosity X-ray sources in the ONC owing to this difference.
To summarize this section, we find that KH 15D is the weakest X-ray emitter known for stars of its spectral class in NGC 2264 and lies about 1 order of magnitude below the median value for the cluster. It is also a weaker X-ray source than all but 3 of the 74 known mid-K star members of the ONC, and two of them have visual extinctions exceeding 5 magnitudes! Again, it lies about an order of magnitude below the median of the cluster. It is possible, of course, that the Chandra exposure we analyzed was obtained, by chance, at a time when KH 15D was at or near the bottom of its range of X-ray variability. It is believed that much of the scatter seen in PMS X-ray luminosity is caused by actual time variations associated with flaring (Feigelson & Montmerle, 1999; Feigelson et al., 2003). This would be an extraordinarily unlucky circumstance, of course, since it is such an extreme outlier, and we consider instead, in the remainder of the paper, whether some aspect of the properties of this star which makes it unique in other ways could also account for its unusually low X-ray luminosity.
## 3 Discussion
Can the low X-ray luminosity of KH 15D be attributed to some sort of extinction effect, perhaps within the circumbinary disk which surrounds it? We find this difficult to support for two reasons. First, there is essentially no reddening or obscuration evident in the light of the K6–7 star during maximum brightness, when the X-ray data were obtained. The star has a color excess of $`E(BV)=0.1`$ mag if it is a K6 star and less if it is K7. This is consistent with what is found for other members of NGC 2264, in which the reddening is known to be small (Rebull et al., 2002; Lamm et al., 2004). If there is any local extinction associated with circumstellar matter, it must be very small. Out of eclipse the star also shows very small photometric variations (less that 0.1 mag in $`I`$) and no detectable color variations (Hamilton et al., 2005). Also, there is no evidence in the X-ray data for a deficiency of soft X-rays which would be most susceptible to absorption. The hard-to-soft ratio is typical of what is found for lightly reddened T Tauri stars such as those in the ONC. We conclude that KH 15D is almost certainly an intrinsically weak X-ray source because of an anomalously low production rate, not because of absorption.
Since the cause of X-ray emission in PMS stars is not fully established, it is not immediately evident how to interpret the low luminosity of KH 15D. Here we discuss two possible explanations, neither of which is without difficulty. It seems likely that, in some way, the binary nature of the star is an important element, so we begin there. An attractive unifying paradigm for the unique and, in some cases, anomalous properties of this WTTS is provided by the eccentric binary model of Winn et al. (2004). It is shown by these authors that constraints on the system from the historic and current light curves can be understood if KH 15D is a roughly equal luminosity binary system with a highly eccentric ($`e0.50.8`$) orbit and period of 48.4 days. The orbit is, at present, slightly inclined to the plane of a circumbinary disk so that one (and only one) of the stars periodically rises above it. Precession of the disk plane is plausibly responsible for the secular variation in the eclipse duration. This compelling model implies a separation of the components at periastron of only about 0.08 AU, close enough to consider possible tidal effects or other influences that such a close approach could have on the system and, in particular, X-ray production.
### 3.1 Interacting Magnetospheres at Periastron?
The radius of the only currently visible star is about 1.3 $`R_{}`$ based on its luminosity and effective temperature (Hamilton et al., 2005). The currently invisible companion was last seen in 1995 and measured to be brighter by several tenths of a magnitude than the K7 star. The historical light curve of the system also demonstrates that the unseen companion is slightly more luminous than the visible star. Assuming the stars are coeval, which seems inescapable, simple theoretical considerations demand that the unseen star be slightly more massive and larger than the K7 star. Hence its radius is probably a little larger than 1.3 $`R_{}`$ but not much larger. The separation of the two components at periastron is about 15 stellar radii. Since magnetospheres of WTTS are typically believed to extend 5–10 stellar radii from the surfaces (Ostriker & Shu, 1995; Preibisch et al., 2005) disruption of the magnetosphere by interactions with matter (or magnetic fields) inside of this point could play a role in lessening X-ray luminosity either by cooling or by lack of confinement of hot gas.
Unfortunately there is little evidence to support (or refute) this hypothesis in the observed X-ray luminosity of other PMS binaries. Perhaps the most similar system known is DQ Tau, which is a CTTS with nearly equal-mass mid-K components in an orbit of eccentricity e=0.56, which brings the stars within about 8 stellar radii of each other at periastron (Mathieu et al., 1997). The star is not detected in the ROSAT All-Sky Survey (Kon̈ig, Neuhaüser, & Stelzer, 2001), which means it is weaker than many PMS stars in Taurus. However, since it is a CTTS it is possible that its low X-ray luminosity is due to absorption in circumstellar matter. Another reasonably eccentric (e=0.24) PMS binary is UZ Tau E (Martin et al., 2005). Unfortunately, it is only a few arc-seconds from UZ Tau W (also a binary) and so the X-ray luminosity of this binary is not measured. Its separation at periastron is also a little larger (about 25 solar radii), it has a much smaller mass ratio (q=0.2) and it is also a CTTS, so there are potentially important differences with KH 15D.
One WTTS spectroscopic binary with a K7 primary, circular orbit and separation of 12.6 stellar radii, V826 Tau, is observed to be roughly normal in its X-ray luminosity, with a quiescent luminosity of around $`2\times 10^{30}`$ erg s<sup>-1</sup> (Reipurth et al., 1990; Carkner et al., 1996). This shows that, proximity of stars, by itself, may not be sufficient to disrupt X-ray emission. However, it may be the variation of the magnetic influence caused by an eccentric orbit that is the key to disrupting a dynamo, so V826 Tau may also not be the best analogue. Since there is no observational evidence that proximity of magnetospheres is a sufficient cause to reduce X-ray emission, we consider another aspect of close periastron passages, namely tidal interactions and possible rotational synchronization.
### 3.2 Tidally Influenced Rotation?
Rotation is a factor in the X-ray luminosity of stars as young as 30 My but it has not been proven to be important in T Tauri stars; evidence to date suggests it is not. Several authors find no correlation between rotation and X-ray emission for PMS stars in the ONC (Gagne & Caillault, 1994; Feigelson et al., 2003). Perhaps the large and variable extinction effects as well as the difficulty of discriminating between WTTS and CTTS in that cluster cause problems with the interpretation. In this regard it will be interesting to see what studies in slightly older and less highly obscured regions such as the Orion Flanking Fields and NGC 2264 will reveal about the role of rotation in X-ray production (Ramirez et al., 2004). Since X-ray emission in TTS is not yet fully understood and since rotation could be a factor in at least some stars, we inquire whether the rotation of the visible star in the KH 15D system is unusual in any way.
There are two methods for determining the rotation rate of a WTTS and both have been employed in the case of KH 15D. Most directly, one can search for periodic fluctuations in the stellar brightness associated with the rotation of a spotted surface. This has been done by Hamilton et al. (2005) and they have detected two clearly significant peaks in the periodogram of out-of-eclipse data at two separate epochs. In both cases, the period was 9.6 days, strongly suggesting that this is, in fact, the rotation period of the visible component of the binary. Confirmation of that comes from a new measurement of v sin i, based on high resolution spectra taken out of eclipse at the Keck and McDonald observatories by the same group. Hamilton et al. (2005) find a value of $`v`$ sin $`i`$ = 6.9 $`\pm `$ 0.3 km s<sup>-1</sup> (replacing an earlier estimate of v sin i $`<`$ 5 km s<sup>-1</sup> by Hamilton et al. (2003) that did not take proper account of macroscopic turbulence in the comparison star). Combining the new v sin i measurement with the known radius of the star, $`R=1.3R_{}`$, yields an expected rotation period ($`P`$) of $`P`$ sin $`i`$ = $`9.4\pm 0.3`$ days (Hamilton et al., 2001, 2003). Since sin $`i1`$ for this eclipsing system, we conclude that KH 15D has a rotation period of 9.6 $`\pm `$ 0.1 days.
This rotation period of KH 15D is somewhat long for a WTTS of its mass in NGC 2264, where the (bi)modal values are near 1 and 4 days (Lamm et al., 2004). It is not, however, the slowest rotator in the cluster. Of the 184 stars with R-I $`<`$ 1.84 (roughly corresponding to mass $`>`$ 0.25 M) in the study by Lamm et al. (2005) 22 (12%) have periods of 9.6 days or longer. If no correlation between rotation period and X-ray emission exists among NGC 2264 stars in general then the significance of KH 15D’s slower than usual rotation for the problem discussed here, is obscure. We note, however, that because it is a member of a relatively close binary, the rotation of the visible component may have been affected by tidal interaction with its primary and could be tidally synchronized (or, rather, pseudosynchronization) as discussed by Hamilton et al. (2005).
Hut (1981) has shown that the pseudosynchronzation angular rotation frequency is a nearly constant fraction (f) of the orbital angular frequency at periastron for orbits in the eccentricity range, e = 0.3 to 0.8. One may write, therefore, that
$$P_{ps}=\frac{P_{orb}}{f}\frac{(1e^2)^{\frac{3}{2}}}{(1+e)^2}.$$
Identifying P<sub>ps</sub> as the measured rotation period, P<sub>orb</sub> as the orbital period and taking f=0.81, as appropriate to the plausible eccentricity range of KH 15D (Johnson et al., 2004). we find the equation is satisfied for e = 0.65 $`\pm `$ 0.01. This is slightly outside the range of solutions (e = 0.68 to 0.8) favored by Johnson et al. (2004) based on astrophysical grounds (primarily the system’s total mass and mass ratio), but well within the plausible range based on the radial velocity curve. It is also consistent with the best fit eccentricity, that comes from modeling the historical and modern light curves (Winn et al., 2005). An estimate of the time scale for pseudosynchronization based on the work of Zahn (1977) suggests that this could have been achieved within a couple of My, as would be required by the inferred age of the KH 15D system.
To summarize, we have found that KH 15D is an unusual stellar system in a new way — it is a very weak source of X-ray emission for its mass and age. It seems likely to us that the eccentric binary nature and close periastron approach are probably involved in this. One possible mechanism is disruption of the magnetosphere of the visible star (and probably both stars) during repeated periastron passages due to magnetic reconnection events. Another, perhaps more likely, possibility is disruption of the usual magnetic dynamo through tidal interactions which could also be implicated in the slower than normal rotation of the visible component. Of course, these are not mutually exclusive mechanisms. Further observations are needed to establish the degree to which KH 15D is in fact anomalous in its X-ray properties for a WTTS and whether either proposed mechanism, or perhaps both, can indeed account for the dearth of X-rays emission.
We are deeply indebted to the referee, John Stauffer, for his detailed and enormously helpful report on the original manuscript including providing some of the comparison data for Figure 1. We are likewise indebted to Luisa Rebull for her part in making that comparison sample available to us. We thank Mike Simon, Bob Mathieu, Eric Jensen, Soeren Meibom and Josh Winn for helpful suggestions related to binary PMS stars. We thank Eric Feigelson for helpful discussions and for his leadership of the COUP survey which produced the important ONC comparison sample of Fig. 2. This material is partly based on work supported by the National Aeronautics and Space Administration under Grant NAG5-12502 issued through the Origins of Solar Systems Program.
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# Luminosity-dependent evolution of soft X–ray selected AGN
## 1 Introduction
In recent years the bulk of the extragalactic X–ray background in the 0.1-10 keV band has been resolved into discrete sources with the deepest ROSAT, Chandra and XMM–Newton observations (Hasinger et al. has98 (1998), Mushotzky et al. mus00 (2000), Giacconi et al. gia01 (2001, 2002), Hasinger et al. has01 (2001), Alexander et al. ale03 (2003), Worsley et al. wor04 (2004)). Optical identification programmes with Keck (Schmidt et al. schm98 (1998), Lehmann et al. leh01 (2001), Barger et al. bar01 (2001, 2003)) and VLT (Szokoly et al. szo04 (2004), Fiore et al. fio03 (2003)) find predominantly unobscured AGN–1 at X–ray fluxes $`S_X>10^{14}`$ erg cm<sup>-2</sup> s<sup>-1</sup>, and a mixture of unobscured AGN–1 and obscured AGN–2 at fluxes $`10^{14}>S_X>10^{15.5}`$ erg cm<sup>-2</sup> s<sup>-1</sup> with ever fainter and redder optical counterparts, while at even lower X–ray fluxes a new population of star forming galaxies emerges (Hornschemeier et al. hor00 (2003), Rosati et al. ros02 (2002), Norman et al. nor04 (2004)). At optical magnitudes R$`>`$24 these surveys suffer from large spectroscopic incompleteness, but deep optical/NIR photometry can improve the identification completeness significantly, even for the faintest optical counterparts (Zheng et al. zhe04 (2004), Mainieri et al. mai05 (2005)).
The AGN/QSO luminosity function and its evolution with cosmic time are key observational quantities for understanding the origin of and accretion history onto supermassive black holes, which are now believed to occupy the centers of most galaxies. X–ray surveys are practically the most efficient means of finding active galactic nuclei (AGNs) over a wide range of luminosity and redshift. Enormous efforts have been made by several groups to follow up X–ray sources with major optical telescopes around the globe, so that now we have fairly complete samples of X–ray selected AGNs. In this work we concentrate on unabsorbed (type–1) AGN selected in the soft (0.5–2 keV) X–ray band, where due to the previous ROSAT work (see Miyaji et al., paper1 (2000, 2001), hereafter Paper I and II) complete samples exist, with sensitivity limits varying over five orders of magnitude in flux, and survey solid angles ranging from the whole high galactic latitude sky to the deepest pencil-beam fields. These samples enable us to construct and probe luminosity functions over cosmological timescales, with an unprecedented accuracy and parameter space.
Conceptually, space densities and luminosity functions are simply derived by dividing the observed number of objects $`N`$ by the volume $`V`$, in which they have been surveyed. The binned luminosity function derived in this fashion generally does not represent the center of the bin. In paper II, we introduced the estimator $`N_{\mathrm{obs}}/N_{\mathrm{mdl}}`$, where $`N_{\mathrm{mdl}}`$ is the number of objects expected from an analytical representation of the luminosity function. Scaling the model value of the analytical function at the bin center by the estimator removes the binning bias to first order. This method is applied in Sect. 4. A quite different method based on $`1/V_{\mathrm{max}}`$ values (Schmidt schmidt68 (1968)) for individual objects is used in Sect. 5. It involves a derivation of the zero-redshift luminosity function that is free of binning bias. The luminosity function at higher redshifts is again derived by employing an analytical representation of the luminosity function and scaling it by the ratio of observed over expected numbers. The use of individual $`1/V_{\mathrm{max}}`$ values in this case allows accounting for an effective optical magnitude limit beyond which redshifts have generally not been obtained in some surveys.
Throughout this work we use a Hubble constant $`H_0=70h_{70}`$ $`\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ and cosmological parameters $`(\mathrm{\Omega }_\mathrm{m},\mathrm{\Omega }_\mathrm{\Lambda })=`$ (0.3, 0.7) consistent with the WMAP cosmology (Spergel et al. spe03 (2003)).
## 2 The X–ray selected AGN–1 sample
For the derivation of the X–ray luminosity function and cosmological evolution of AGN we have chosen well–defined flux–limited samples of active galactic nuclei, with flux limits and survey solid angles ranging over five and six orders of magnitude, respectively. To be able to utilize the massive amount of optical identification work performed previously on a large number of shallow to deep ROSAT surveys, we restricted the analysis to samples selected in the 0.5–2 keV band. In addition to the ROSAT surveys already used in Paper I and II, we included data from the recently published ROSAT North Ecliptic Pole Survey (NEPS, Gioia et al. gio03 (2003), Mullis et al. mul04 (2004)), from an XMM–Newton observation of the Lockman Hole (Mainieri et al. mai02 (2002)) and the Chandra Deep Fields South (CDF–S, Szokoly et al. szo04 (2004), Zheng et al. zhe04 (2004), Mainieri et al. mai05 (2005)) and North (CDF–N, Barger et al. bar01 (2001, 2003)). In order to avoid systematic uncertainties introduced by the varying and a priori unknown AGN absorption column densities we selected only unabsorbed (type–1) AGN, classified by optical and/or X–ray methods. We are using here a definition of type–1 AGN, which is largely based on the presence of broad Balmer emission lines and small Balmer decrement in the optical spectrum of the source (optical type–1 AGN, e.g. the ID classes a, b, and partly c in Schmidt et al. schm98 (1998)), which largely overlaps the class of X–ray type–1 AGN defined by their X–ray luminosity and unabsorbed X–ray spectrum (Szokoly et al. szo04 (2004)). However, as Szokoly et al show, at low X–ray luminosities and intermediate redshifts the optical AGN classification often breaks down because of the dilution of the AGN excess light by the stars in the host galaxy (see e.g. Moran et al. (mor02 (2002)), so that only an X–ray classification scheme can be utilized. Schmidt et al. (schm98 (1998)) have already introduced the X–ray luminosity in their classification. For the deep XMM–Newton and Chandra surveys we in addition use the X–ray hardness ratio to discriminate between X–ray type–1 and type–2 AGN, following Szokoly et al. (szo04 (2004)).
In order to convert the count rates observed in the 0.5–2 keV band to unabsorbed 0.5–2 keV fluxes, we assumed a power law AGN spectrum with a photon index of $`\mathrm{\Gamma }=2.0`$ and Galactic absorption. Typical values of the Galactic neutral hydrogen column density are $`(0.51)\times 10^{20}`$ $`\mathrm{cm}^2`$ for the deep surveys and a maximum of $`16\times 10^{20}`$ $`\mathrm{cm}^2`$ for a small portion of the sky covered by the ROSAT Bright Survey (RBS). Because the band, in which the AGN are selected, is the same as the one for which we calculate the fluxes, systematic differences in the true AGN power law indices have a negligible effect on the derived fluxes. Assuming spectral indices in the range $`\mathrm{\Gamma }=13`$, the conversion between the observed 0.5–2 keV count rates and the X–ray flux $`S_\mathrm{x}`$ (here and hereafter, $`S_\mathrm{x}`$ represents the 0.5–2 keV flux and $`S_{\mathrm{X14}}`$ is the same quantity measured in units of $`10^{14}\mathrm{erg}\mathrm{s}^1\mathrm{cm}^2`$) varies by less than 10%.
The surveys we have used are summarized in Table 1. A total of 944 X–ray selected type–1 AGN were compiled from eight independent samples containing a total of 2566 soft X–ray sources. The number of unidentified sources<sup>1</sup><sup>1</sup>1We call unidentified sources those, which do not have a reliable redshift determination, either through spectroscopy or through photometric redshifts; however practically all of these have optical or NIR counterparts in these samples is only 86 (of which 57 could be AGN–1), yielding an unprecedented identification fraction of 97%. Due to the extreme faintness of the optical counterparts, the lowest identification fractions are achieved in the recent deepest samples: 87% for the XMM–Newton survey in the Lockman Hole and 88% in the CDF–N. A surprisingly high identification fraction of 98% has been achieved in the CDF–S through the utilization of photometric redshifts based on extremely faint optical/NIR photometry.
For the computation of the soft X-ray luminosity function SXLF, it is important to define the available survey solid angle as a function of limiting flux. In case there is incompleteness in the spectroscopic identifications in the ROSAT surveys, we have made the usual assumption that the redshift/classification distribution of these unidentified sources is the same as the identified sources at similar fluxes by defining the ’effective’ survey solid angle as the geometrical survey solid angle multiplied by the completeness of the identifications (see Paper I). This assumption is not correct when the source is unidentified due to non–random causes, in particular its optical faintness. The treatment of this identification incompleteness and the effect of the optical limit on the derived space densities is discussed in detail in Sect 5.
In several surveys we had to choose an X–ray flux limit a posteriori, based on optical completeness criteria, i.e. maximising the number of optically identified sources, while simultaneously minimising the number of unidentified objects. This procedure can introduce a bias against optically faint sources, if the reason for the missing redshift is the optical faintness of the source and in fields with relatively few objects. We have tried to minimise the impact of this ”gerrymandering” effect, e.g. by allowing a number of unidentified sources to enter the sample and then defining the corresponding X–ray flux limits in the geometric mean between the last identified and the next unidentified source. In addition, the wide range for X–ray and thus optical flux limits in our survey tends to reduce biases, which occur at the flux limit of individual surveys. The problem of missing redshifts in the faintest surveys is addressed specifically in Sect 4 and 5.
Below we summarize our sample selection and completeness for each survey. Figure 1 shows the AGN–1 sample in the redshift – luminosity plane. Figure 2 gives the combined solid angle versus flux curve. Both the sample and the solid angle coverage are available in computer readable form under http://mpe.mpg.de/~ghasinger.
### 2.1 The ROSAT Bright Survey (RBS)
The RBS identified the brightest $``$ 2000 X–ray sources detected in the ROSAT All–Sky Survey (RASS, Voges et al. vog99 (1999)) at high galactic latitude, $`|b|>30\mathrm{°}`$, excluding the Magellanic Clouds and the Virgo cluster, with ROSAT PSPC count rates above 0.2 s<sup>-1</sup>. This program achieved a spectroscopic completeness of 99.5% (Schwope et al. schw00 (2000)). We selected the sub–sample of 931 sources with count rates above 0.2 s<sup>-1</sup> in the ROSAT 0.5–2 keV band (PSPC channels 52-201), which is 100% identified. Since the absorption in our galaxy varies from place to place, the same count rate limit corresponds to different 0.5–2 keV flux limits based on the different galactic $`N_\mathrm{H}`$ values. The $`N_\mathrm{H}`$ value ranges from $`(0.516)\times 10^{20}`$ $`\mathrm{cm}^2`$ in the RBS survey area. Correspondingly, the survey solid angle varies steeply with flux from about 3000 deg<sup>2</sup> at a flux limit of $`S_{\mathrm{X14}}=246`$ to a total of 20391 deg<sup>2</sup> at a flux limit of $`S_{\mathrm{X14}}=360`$.
### 2.2 The RASS Selected-Area Survey North (SA–N)
This survey gives optical identifications of a representative sample of northern ($`\delta >9\mathrm{°}`$) RASS sources in six study areas outside the Galactic plane ($`|b|>19.6\mathrm{°}`$) with a total of $`685\mathrm{deg}^2`$. A count rate limited complete RASS subsample comprising 674 sources has been identified (Appenzeller et al. app98 (1998)). The fields selected for the survey have a Galactic column density in the range $`N_\mathrm{H}=(211)\times 10^{20}`$ $`\mathrm{cm}^2`$. We have further selected our sample such that each of the six fields has a complete ROSAT hard-band (0.5–2 keV; channels 52-201) countrate-limited sample with complete identifications ($`CR_{0.52\mathrm{k}\mathrm{e}\mathrm{V}}>`$ 0.01–0.05 $`\mathrm{cts}\mathrm{s}^1`$). To avoid overlap with the RBS, those sources common in both samples were removed from SA–N, yielding a total of 406 sources with 98.5% spectroscopic completeness.
### 2.3 The ROSAT North Ecliptic Pole Survey (NEPS)
The RASS data in a contiguous area of 80.7 deg<sup>2</sup> around the North Ecliptic Pole (Galactic latitude $`b>29.8\mathrm{°}`$) have been used to construct a survey consisting of 445 X–ray sources detected above a 4$`\sigma `$ threshold. Gioia et al. (gio03 (2003)) and Mullis et al. (mul04 (2004)) have identified 99.6% of these sources and determined redshifts for the extragalactic objects. Since the exposure in the ROSAT All-Sky Survey increases significantly towards the North-Ecliptic Pole, the actual survey sensitivity is a strong function of ecliptic latitude. The original NEPS sample is selected in the full ROSAT PSPC band. For consistency with the other surveys used in our work, we selected sources detected significantly in the PSPC hard band (0.5–2 keV; channels 52-201) by specifying a hard count rate limit as a function of ecliptic latitude. New 0.5–2 keV fluxes were calculated from the hard PSPC count rates in the same way as for the RBS and SA–N objects, taking into account the Galactic neutral hydrogen column density varying in the range $`2.66.2\times 10^{20}cm^2`$ across the NEPS region. Due to the large gradient in exposure times, we have cut the sample to 252 sources with fluxes above $`S_{\mathrm{X14}}=10.1`$, where the solid angle of this survey is 70.5 deg<sup>2</sup>, increasing to 80.8 deg<sup>2</sup> at $`S_{\mathrm{X14}}=12.4`$. Only one of these sources remains unidentified in the NEPS sample.
### 2.4 The ROSAT International X–ray/Optical Survey (RIXOS)
The ROSAT International X–ray/Optical Survey (RIXOS, Mason et al. mas00 (2000)) is a medium-sensitivity survey and optical identification program of X–ray sources discovered in ROSAT high Galactic latitude fields ($`|b|>28\mathrm{deg}`$) and observed with the Position Sensitive Proportional Counter (PSPC) detector with a minimum exposure time of 8 ks. The survey comprises 82 ROSAT PSPC fields and made use of the central 17 arcmin of each field, however, excluding the target region for pointings on known X–ray sources. The total survey contains 395 X–ray sources, selected in the PSPC 0.5–2 keV band. A flux limit of $`S_{\mathrm{X14}}=3.0`$ was adopted for the survey, substantially above the detection threshold of each field, however, the actual spectroscopic completeness limit varies from field to field. We have chosen a strategy, which on one hand maximises the sample of identified AGN–1 and on the other hand minimises the number of unidentified sources. There are 51 fields (12.3 deg<sup>2</sup>) identified completely down to the survey flux limit. Three fields have such a low identification fraction, that we ignore them. For the remaining 28 fields we allow at most one unidentified source. If an unidentified source has the lowest flux of the subsample of a particular field, we exclude this source and raise the flux limit for this field to the geometric average between the flux of this source and that of the last identified source. This way we can define a clean RIXOS sample comprising 340 objects and only 14 unidentified sources, i.e. an identification fraction of 95.9%. The survey solid angle, corrected for spectroscopic incompleteness, rises from 15.0 deg<sup>2</sup> at $`S_{\mathrm{X14}}=3.0`$ to 19.5 deg<sup>2</sup> at $`S_{\mathrm{X14}}=10.2`$.
### 2.5 The ROSAT Medium Survey (RMS)
For this work we have grouped a number of medium–deep ROSAT surveys with flux limits in the range $`S_{\mathrm{X14}}=0.51`$ into the RMS. In particular these comprise pointed observations at the North Ecliptic Pole (Bower et al. bow96 (1996)), the UK Deep Survey (McHardy et al. mch98 (1998)), the Marano field (Zamorani et al. zam99 (1999)) and the outer parts of the Lockman Hole (Schmidt et al. schm98 (1998), Lehmann et al. leh00 (2000)). The North Ecliptic pole pointing covers the same sky area as the center of the NEPS, however, to a flux limit of $`S_{\mathrm{X14}}=1.0`$. Again, we remove the overlapping sources between the two surveys. For the UK Deep Survey and the Marano Field we define a flux limit of $`S_{\mathrm{X14}}=0.5`$, following Paper I. For the ROSAT PSPC survey of the Lockman Hole we only chose the region not covered by the deeper RDS/XMM survey (see section 2.6) but otherwise selected the completely identified sample with the same flux limits as those chosen for the ROSAT Ultradeep Survey UDS by Lehmann et al. leh01 (2001): $`S_{\mathrm{X14}}>0.96`$ for PSPC off–axis angles in the range 12.5–18.5 arcmin and $`S_{\mathrm{X14}}>0.55`$ for off–axis angles smaller than 12.5 arcmin. Overall, the RMS contains 124 sources, at an identification completeness of 94.4%. Correspondingly, the corrected survey solid angle varies in the range 0.30–0.70 deg<sup>2</sup> for flux limits $`S_{\mathrm{X14}}`$=0.5–1.
### 2.6 Deep XMM–Newton survey of the Lockman Hole (XMM/RDS)
The Lockmam Hole (XMM/RDS) has been observed by XMM–Newton a total of 17 times during the PV, AO–1 and AO–2 phases of the mission, with total good exposure times in the range 680–880 ks in the PN and MOS instruments (see Hasinger et al. has01 (2001, 2004) and Worsley et al. wor04 (2004) for details). Spectroscopic optical identifications of the ROSAT sources in the LH have been presented by Schmidt et al. (schm98 (1998)) and Lehmann et al. (leh00 (2000, 2001)) and a new catalogue from the XMM–Newton PV phase is given in Mainieri et al. (mai02 (2002)). Some photometric redshifts have been discussed in Fadda et al., (fad02 (2002)). Here we selected sources from the 770 ksec dataset (Brunner et al., 2005, in prep.) with additional spectroscopic identifications obtained with the DEIMOS spectrograph on the Keck telescope in spring 2003 and 2004 by M. Schmidt and P. Henry (Szokoly al., 2005, in prep.). In order to maximise the spectroscopic/photometric completeness of the sample, we selected objects in two off–axis intervals: $`S_{\mathrm{X14}}=0.38`$ for off–axis angles in the range 10.0–12.5 arcmin and $`S_{\mathrm{X14}}=0.13`$ for off–axis angles smaller than 10.0 arcmin. The total number of sources in the XMM/RDS is 81, with 8 potential AGN–1 still unidentified.
### 2.7 Chandra Deep Field South (CDF–S)
We have used the catalogue of Giacconi et al. (gia02 (2002)) based on the 1 Ms observation of the CDF–S (Rosati et al. ros02 (2002)). Spectroscopic identifications with the FORS instruments at the ESO VLT have been obtained by Szokoly et al. (szo04 (2004)), yielding a spectroscopic completeness around 60%. Additional spectroscopic redshifts of CDF–S X–ray sources have been obtained with the VIMOS spectrograph at the ESO VLT (Lefevre et al., 2004). The field is also included in the COMBO-17 intermediate–band optical survey, which gives very reliable photometric redshifts for the brighter sources (Wolf et al., combo17 (2004)). Very deep NIR photometry has been obtained with the ISAAC camera at the VLT in conjunction with deep optical imaging with the HST ACS as part of the GOODS project (Dickinson & Giavalisco dic03 (2003), Mobasher et al. mob04 (2004)). The CDF–S therefore offers the highest quality photometric redshifts of faint X–ray sources, which are discussed in Zheng et al. (zhe04 (2004)) and Mainieri et al. (mai05 (2005)). Using the ISAAC images, tentative photometric redshifts could even be assigned to several of the extreme X–ray/optical sources (EXOs) discussed by Koekemoer et al. (koe04 (2004)). We selected all sources from the Giacconi et al. gia02 (2002) catalogue within 10 arcmin from the Chandra pointing center significantly detected in the 0.5–2 keV band. The sample thus contains a total of 293 objects. Combining all spectroscopic and photometric redshifts, only 2 sources in the CDFS remain unidentified, of which one could be an AGN–1. The survey solid angle for the CDF–S has been estimated using a simple off–axis dependent flux limit. The solid angle increases from 0.023 deg<sup>2</sup> at $`S_{\mathrm{X14}}=0.0053`$ to 0.087 deg<sup>2</sup> at $`S_{\mathrm{X14}}=0.027`$.
### 2.8 Chandra Deep Field North (CDF–N)
We have used selected X–ray sources from the 2 Ms CDF–N source catalogue by Alexander et al. (ale03 (2003)) along with optical identifications by Barger et al. (bar03 (2003)) for our AGN–1 sample. Following Szokoly et al. (szo04 (2004)), we selected AGN–1 either from broad permitted Balmer lines or from the X–ray luminosity and hardness, using HR$`<`$–0.2. We set our flux limits such that we also have sufficient hard (2–8 keV) sensitivity to exclude objects with HR$`>`$–0.2 and to include as many sources as possible which meet these criteria. Unlike the CDF–S, the CDF–N exposure map has a complicated structure and a simple off–axis dependence of the limiting flux is not a good approximation. Thus we have used the rectangular region of 170 arcmin<sup>2</sup>, which has the deepest coverage and is mostly co–spatial with the region covered by HST ACS in the GOODS project (Giavalisco et al. goods1 (2004), Cowie et al. cowie04 (2004)). Within this region, the exposure and background are smooth enough that the photon counts limit of the detected sources can be approximated by a simple function of off–axis angle. In practice, due to statistical fluctuations, three sources have upper limits to HR between –0.1 and –0.2 and those have been considered to meet our hardness ratio criterion. Among the 128 sources meeting the soft counts limit and hardness ratio criteria, 20 are unidentified and 5 are stars (85% completeness). Only one broad–line AGN had a harder hardness ratio than our limit; this was also included in our type–1 AGN sample. The flux–solid angle relation has been calculated from the “limiting flux map”, where the counts limit is divided by the soft–band exposure map (in seconds) and multiplied by the conversion factor of $`510^{12}\mathrm{erg}\mathrm{s}^1\mathrm{cm}^2`$ (Alexander et al., ale03 (2003)). Due to the incompleteness in this field, where most unidentified sources are optically faint, this sample has not been included in the analysis in Section 4, but considered in Section 5, where a method is developed to account for the optical magnitude limit in calculating the survey volume.
## 3 Number counts for different source classes
The combination of a large number of surveys with a wide range of sensitivity limits and solid angle coverage presents a unique resource. On one hand, the surveys presented here resolve the soft X–ray background almost completely. On the other hand, we have an almost complete optical identification and redshift determination for all constituents. For the first time in any astrophysical waveband we are thus in the position to study the complete contribution of different object classes to the X–ray background and their evolution with cosmic time. Using the solid angle versus flux limit curve given in Figure 2 we compiled number counts for the total sample including all classes of sources and for the subclass of AGN–1. Figure 3 (top) shows the cumulative source counts. For clarity we also show normalized differential source counts $`dN/dS_{\mathrm{X14}}S_{\mathrm{X14}}^{2.5}`$ in the bottom panel of Figure 3. Euclidean source counts would correspond to horizontal lines in this graph.
For the total source counts, the well-known broken power law behaviour is confirmed with high precision. We fitted a broken power law to the differential source counts and obtain power law indices of $`\alpha _b=2.34\pm 0.01`$ and $`\alpha _f=1.55\pm 0.04`$ for the bright and faint end, respectively, a break flux of $`S_{\mathrm{X14}}=0.65\pm 0.10`$ and a normalisation of dN/d$`S_{\mathrm{X14}}=103.5\pm 5.3`$ deg<sup>-2</sup> at $`S_{\mathrm{X14}}=1.0`$ with a reduced $`\chi ^2`$=1.51. The total differential source counts, normalized to a Euclidean behaviour (dN/d$`S_{\mathrm{X14}}\times S_{\mathrm{X14}}^{2.5}`$ is shown with open symbols in Figure 3. We see that the total source counts at bright fluxes, as determined by the ROSAT All-Sky Survey data, are significantly flatter than Euclidean, consistent with the discussion in Hasinger et al. (has93 (1993)). Moretti et al. (mor03 (2003)), on the other hand, have derived a significantly steeper bright flux slope ($`\alpha _b2.8`$) from ROSAT HRI pointed observations. This discrepancy can probably be attributed to the selection bias against bright sources, when using pointed observations where the target area has to be excised.
Type–1 AGN are the most abundant population of soft X–ray sources. For the determination of the AGN–1 number counts we include those unidentified sources, which have hardness ratios consistent with AGN–1 (a contribution of $`6\%`$, see Table 1). Figure 3 shows, that the break in the total source counts at intermediate fluxes is produced by type–1 AGN, which are the dominant population there. Both at bright fluxes and at the faintest fluxes, type–1 AGN contribute about 30% of the X–ray source population. At bright fluxes, they have to share with clusters, stars and BL-Lac objects, at faint fluxes they compete with type–2 AGN and normal galaxies. We fitted a broken power law to the differential AGN–1 source counts and obtain power law indices of $`\alpha _b=2.55\pm 0.02`$ and $`\alpha _f=1.15\pm 0.05`$ for the bright and faint end, respectively, a break flux of $`S_{\mathrm{X14}}=0.53\pm 0.05`$, consistent with that of the total source counts within errors, and a normalisation of of dN/d$`S_{\mathrm{X14}}=83.2\pm 5.5`$ deg<sup>-2</sup> at $`S_{\mathrm{X14}}=1.0`$ with a reduced $`\chi ^2`$=1.26. The AGN–1 differential source counts, normalized to a Euclidean behaviour (dN/d$`S_{\mathrm{X14}}\times S_{\mathrm{X14}}^{2.5}`$) is shown with filled symbols in Figure 3. Also shown are the predictions of the best–fit SXLF models discussed in Section 5.
## 4 The SXLF and the Space Density Function
### 4.1 Basic method
In this section, we present the binned Soft X–ray Luminosity Function (SXLF) of type–1 AGNs. The basic approach is to use the $`N^{\mathrm{obs}}/N^{\mathrm{mdl}}`$ estimator described in Paper II. The procedure is outlined below:
1. Divide the combined sample into several redshift shells. For each redshift shell, fit the AGN XLF with a smooth analytical function using a Maximum-likelihood fit over each object (i.e., without binning; see Paper I for details).
2. For the fitted model in each redshift shell, check the absolute goodness of fit with one– and two–dimensional Kolmogorov–Smirnov tests (hereafter, 1D–KS and 2D–KS tests respectively; Press et al. numrec (1992), Fassano & Franceschini ff\_2dks (1987)). The K–S tests are also for the unbinned data sets and thus are free from artefacts and biases due to binning.
3. For each redshift shell, bin the objects in luminosity bins to determine the observed number of objects ($`N^{\mathrm{obs}}`$).
4. For each luminosity bin, evaluate the analytical fit at the central luminosity/redshift ($`d\mathrm{\Phi }^{\mathrm{mdl}}/d\mathrm{log}L_\mathrm{x}`$).
5. Calculate the predicted number of AGNs in the bin ($`N^{\mathrm{mdl}}`$).
6. The final result is
$$d\mathrm{\Phi }/d\mathrm{log}L_\mathrm{x}=d\mathrm{\Phi }^{\mathrm{mdl}}/d\mathrm{log}L_\mathrm{x}N^{\mathrm{obs}}/N^{\mathrm{mdl}}$$
(1)
For the analytical expression of the SXLF in each redshift shell, we use the smoothed two power law formula. Because the redshift shells have a finite widths, the fit results depend on the evolution of the SXLF within them:
$$\frac{\mathrm{d}\mathrm{\Phi }(L_\mathrm{x},z)}{\mathrm{d}\mathrm{log}L_\mathrm{x}}\left[\left(\frac{L_\mathrm{x}}{L_{\mathrm{x},}}\right)^{\gamma _1}+\left(\frac{L_\mathrm{x}}{L_{\mathrm{x},}}\right)^{\gamma _2}\right]^1e_\mathrm{d}(z,L_\mathrm{x}),$$
(2)
where $`e_\mathrm{d}(z,L_\mathrm{x})`$ is the density evolution factor. While the final results are insensitive to the detailed behavior of $`e_\mathrm{d}(z,L_\mathrm{x})`$ within the shell at most locations in the $`(L_\mathrm{x},z)`$ space, we have taken our best-estimate by using the luminosity-dependent density evolution (LDDE) model derived later in Sect. 4.4. The luminosity range of the fit is from log$`L_\mathrm{x}`$=42.0 to the maximum available luminosity in the sample.
In this section, we tried to make the sample as complete as possible, and we excluded the CDF–N from the analysis, where the incompleteness fraction is significant and most of the unidentified sources are optically faint. All of the unidentified sources in the ROSAT samples are optically bright and the reasons for them to be unidentified are mostly by random causes, i.e., are not correlated with the intrinsic properties of the source. For the CDF–S, extensive photometric redshift studies including COMBO-17 (Wolf et al., combo17 (2004)) and a careful individual photometric redshift determination of X–ray sources by Zheng et al. zhe04 (2004) and Mainieri et al., mai05 (2005) has left only one potential AGN–1 without redshift information. For RDS/XMM, 2 of the 8 unidentified sources which could be type–1 AGNs from X–ray hardness/spectra criteria are optically bright and they have remained unidentified so far for random reasons. The remaining 6 are optically faint ($`R24.0`$) and the reason for remaining unidentified may well be correlated with redshift.
As our nominal case, we took the first–order approach and defined “effective” survey solid angle (as a function of flux), which is the geometrical survey solid angle multiplied by the completeness, i.e. the fraction of identified X–ray sources in the survey, whether or not they are optically faint or bright. The correction has been made in each survey. In addition, we have also considered the upper bounds on the binned SXLF and the space density function from the sample where all the unidentified optically faint ($`R24.0`$) sources in turn are assigned the central redshift of each redshift shell.
In the latter case, we used the geometrical solid angle for the CDF–S and the incompleteness correction to the RDS/XMM solid angle was only for the optically bright $`R<24`$ unidentified sources.
### 4.2 The binned SXLF
The best fit parameters of Eq. 2 for each redshift shell are shown in Table 2 along with the results of the 1D– and 2D–KS tests (see the notes of the table). The normalization is defined by:
$$A_{44}=\frac{\mathrm{d}\mathrm{\Phi }(L_\mathrm{x}=10^{44}\mathrm{erg}\mathrm{s}^1,z=z_\mathrm{c})}{\mathrm{d}\mathrm{log}L_\mathrm{x}},$$
(3)
where $`z_\mathrm{c}`$ is the central redshift of the shell. The parameter errors in Table 2 correspond to a likelihood change of 2.7 (90% confidence errors), except for the normalization $`A_{44}`$, which cannot be a fit parameter in the maximum-likelihood method. The errors of $`A_{44}`$ are simply taken as the 90% Poisson errors of the number of the sources. Defining the normalization at a fixed luminosity ($`\mathrm{log}L_\mathrm{x}=44`$) minimizes its dependence on other parameters. In any case, Eq. 2 gives a statistically satisfactory expression for all redshift shells as shown in Table 2.
We have made luminosity bins starting with a minimum luminosity of $`\mathrm{log}L_\mathrm{x}=42.0`$ with a smallest bin size of $`\mathrm{\Delta }\mathrm{log}L_\mathrm{x}=0.25`$ in each redshift shell. If there are fewer than 10 AGNs in a bin, we have further rebinned up to a maximum bin size of $`\mathrm{\Delta }\mathrm{log}L_\mathrm{x}=1.0`$. Table 3 shows the full binned results for the nominal case, along with observed number of AGNs ($`N^{\mathrm{obs}}`$), model (Table 2) predicted number ($`N^{\mathrm{mdl}}`$) and final estimated values SXLF value at the center of each bin. For reference, the additional number $`N^f`$ of AGNs for the case where all the optically-faint unidentified sources are assigned the central redshift of the bin (in duplicate, as described above) are also shown in the last column of Table 3. The full SXLF in the 6 redshift shells is plotted in Figs. 4 in separate panels. In all but the closest redshift shell panel, the best–fit two power law function to the $`0.015<z<0.2`$ SXLF (Table 2) are also overplotted for reference. Three overall analytical expressions discussed in Sect. 4.4 are also overplotted for comparison as discussed there. Because of the high completeness of our sample, the redshift distribution of the optically faint sources affects the final SXLF results very little except for the case where all of them happen to fall into the highest redshift shell. In this case, the SXLF in the $`3.2<z<4.8`$, $`44<\mathrm{log}L_\mathrm{x}`$ bin almost double. This is also verified by a comparison with the alternative approach outlined in Sect. 5.
### 4.3 Evolution of the Space Density
In this section, we investigate the evolution of the type–1 AGN space density in different luminosity classes as a function of redshift. The estimator of the space density is the $`N^{\mathrm{obs}}/N^{\mathrm{mdl}}`$. The fit with Eq. 2 has been made in finer redshift shells than in Sect. 4.2. The space densities as a function of redshift were calculated in five luminosity classes with $`\mathrm{log}L_\mathrm{x}`$ of 42–43, 43–44, 44–45, 45–46, and $`>`$46 as well as the sum over all luminosities with $`\mathrm{log}L_\mathrm{x}>42`$. The resulting curves are shown in Fig. 5(a) for the nominal calculations. The incompleteness upper bounds have also been calculated, but have not been shown here for the visibility of the figure. These upper bounds are shown in Sect. 6 (Figs. 11 & 13). Since the Black Hole growth function is more closely linked to the emissivity per comoving volume, we also show the emissivity as a function of redshift in the same luminosity classes in Fig. 5(b).
Figure 5(a) clearly shows a shift of the number density peak with luminosity, in the sense that more luminous AGNs (QSOs) peak earlier in the history of the universe, while the low luminosity ones arise later. Also, there is a clear decline of the derived space densities at least for luminosities of $`\mathrm{log}L_\mathrm{x}<44`$, even when the optical incompleteness upper bounds are taken into account. The counting statistics and spectroscopic incompleteness for the more luminous AGNs do not allow to determine a decline, but do also not exclude it. This issue is further discussed in Sections 5 and 6.
In order to show the behaviour of the luminosity dependence of the evolution more quantitatively, we have also made a maximum-likelihood fit of the evolution curve in each of the luminosity bins, with $`\mathrm{log}L_\mathrm{x}`$ ranges of 42–43,43–44,44–45, and 45–46. We used two power law components of the $`(1+z)`$ evolution with a cutoff redshift:
$$e_\mathrm{d}(z,L_{\mathrm{xc}})=\{\begin{array}{cc}(1+z)^{p1}\hfill & (zz_\mathrm{c})\hfill \\ e_\mathrm{d}(z_\mathrm{c})[(1+z)/(1+z_\mathrm{c})]^{p2}\hfill & (z>z_\mathrm{c})\hfill \end{array}.$$
(4)
where $`L_{\mathrm{xc}}`$ is the central (logarithmic) luminosity of the bin. As was the case for Eq. 2, the fit depends on the shape of the luminosity function (along the luminosity direction) within the luminosity bin. Again, we have fixed the behavior in the luminosity direction using those from the LDDE model (Sect. 4.4) as a template. We also show the normalization:
$$A_0\mathrm{d}\mathrm{\Phi }(L_\mathrm{x}=L_{\mathrm{xc}},z=0)/\mathrm{d}\mathrm{log}L_\mathrm{x}.$$
(5)
The best-fit results are shown in Table 4 together with the K–S probabilities. In the $`\mathrm{log}L_\mathrm{x}=4546`$ bin, the fit for $`p_2`$ was unconstrained. Thus we have fixed the values of $`p_2`$ to that from LDDE (see below) for this luminosity bin.
### 4.4 Global Representations by Analytical Functions
It is sometimes useful to provide a simple analytical fit for the SXLF over the whole redshift-luminosity range. We first used the pure-luminosity evolution (PLE) form, in order to enable a comparison with previous work:
$$\frac{\mathrm{d}\mathrm{\Phi }(L_\mathrm{x},z)}{\mathrm{d}\mathrm{log}L_\mathrm{x}}=\frac{\mathrm{d}\mathrm{\Phi }(L_\mathrm{x}/e_\mathrm{l}(z),0)}{\mathrm{d}\mathrm{log}L_\mathrm{x}},$$
(6)
with the luminosity evolution factor:
$$e_\mathrm{l}(z)=\{\begin{array}{cc}(1+z)^{p1}\hfill & (zz_\mathrm{c})\hfill \\ e(z_\mathrm{c})[(1+z)/(1+z_\mathrm{c})]^{p2}\hfill & (z>z_\mathrm{c})\hfill \end{array}.$$
(7)
We again used the smoothed two power law form (Eq. 2, excluding the $`z`$–dependent factor) for the $`z=0`$ SXLF. The best–fit PLE parameters are shown in Table 5 with the results of the K–S tests. The best–fit PLE model is overplotted with the binned SXLF in Figs. 4 & 5 as dotted lines. It is apparent from the comparison in these figures, especially the latter, PLE does not represent the behaviour of the low–luminosity ($`\mathrm{log}L_\mathrm{x}44`$), intermediate redshift ($`0.5z1.8`$) regime, due to the rather restrictive nature of the PLE form.
As a more general analytical form for a refined representation of the SXLF, we have explored the luminosity-dependent density evolution form (LDDE) form, originally suggested by Schmidt & Green (sg83 (1983)) for describing optically-selected QSOs:
$$\frac{\mathrm{d}\mathrm{\Phi }(L_\mathrm{x},z)}{\mathrm{d}\mathrm{log}L_\mathrm{x}}=\frac{\mathrm{d}\mathrm{\Phi }(L_\mathrm{x},0)}{\mathrm{d}\mathrm{log}L_\mathrm{x}}e_\mathrm{d}(z,L_\mathrm{x}),$$
(8)
where $`e_\mathrm{d}(z,L_\mathrm{x})`$ is the density function normalized to z=0. The results from Sect. 4.3 show that the peak number density shifts from $`z0.7`$ at $`\mathrm{log}L_\mathrm{x}42.5`$ to $`z>2`$ at high luminosities. Based on a similar observation of a hard X–ray–selected sample, Ueda et al. ued03 (2003) used an expression where $`z_\mathrm{c}`$ is a simple function of $`L_\mathrm{x}`$:
$$e_\mathrm{d}(z,L_\mathrm{x})=\{\begin{array}{cc}(1+z)^{p1}\hfill & (zz_\mathrm{c})\hfill \\ e_\mathrm{d}(z_\mathrm{c})[(1+z)/(1+z_\mathrm{c})]^{p2}\hfill & (z>z_\mathrm{c})\hfill \end{array}.$$
(9)
along with
$$z_\mathrm{c}(L_\mathrm{x})=\{\begin{array}{cc}z_{\mathrm{c},0}(L_\mathrm{x}/L_{\mathrm{x},\mathrm{c}})^\alpha \hfill & (L_\mathrm{x}L_{\mathrm{x},\mathrm{c}})\hfill \\ z_{\mathrm{c},0}\hfill & (L_\mathrm{x}>L_{\mathrm{x},\mathrm{c}})\hfill \end{array}.$$
(10)
The results of the analysis in the previous section shown in Table 4 suggest that considering the dependence of $`p1`$ and $`p2`$ on luminosity would still improve the fit. Thus we have also included the following for our full LDDE expression:
$`p1(L_\mathrm{x})=p1_{44}+\beta _1(\mathrm{log}L_\mathrm{x}44)`$ (11)
$`p2(L_\mathrm{x})=p2_{44}+\beta _2(\mathrm{log}L_\mathrm{x}44)`$ (12)
The best–fit parameters and the results of the K–S tests for the PLE and LDDE models are summarized in Table 5. The best–fit PLE and LDDE models are overplotted on Figs. 4 and 5 with dotted and dashed lines respectively. A detailed discussion of the comparison of model and data is given in Sect. 6.
## 5 An alternate approach using the $`V_{\mathrm{max}}`$ method
As described in the Introduction, the luminosity function derived from survey data binned in luminosity and redshift does not necessarily apply to the centers of the $`(L_\mathrm{x},z)`$ bins. This binning bias tends to be especially a problem if data are scarce (often at higher redshifts) and gradients across bins are large. The previous section describes a procedure that corrects the binned space densities to first order.
In this section, we avoid deriving densities from binned survey data. Instead, we use the $`V_{\mathrm{max}}`$ values of individual RBS sources to derive the zero redshift luminosity function. We then derive by iteration an analytical density template at various $`L_\mathrm{x}`$ values that, together with the zero redshift luminosity function, accounts for the observed counts and redshifts of the deeper surveys. The end result of the procedure is a set of observed values of the luminosity function that apply to the centers of the $`(L_\mathrm{x},z)`$ bins, and that is quite insensitive to the precise template employed. A further advantage of employing $`V_{\mathrm{max}}`$ of individual sources is that it can be derived for two or more selection variables. This allows us to account for the effect of a spectroscopic magnitude limit in some of the deeper surveys beyond which the redshift is unknown for most of the sources. In the first use of $`V_{\mathrm{max}}`$, this feature was used to derive the luminosity function of radio quasars from a sample in which only the optically brightest objects had redshifts (Schmidt schmidt68 (1968)).
### 5.1 Using $`V_{\mathrm{max}}`$ to derive the luminosity function
The derivation of a luminosity function from objects in a well defined sample usually involves binning the observations in redshift and luminosity. If we make the bins in luminosity so small that each contains only one or zero objects then the luminosity function is composed of contributions from each of the individual sample objects. In the limit, each of the objects contributes to the luminosity function a delta function of amplitude $`1/V_{\mathrm{max}}`$ at the object’s luminosity $`L`$, where $`V_{\mathrm{max}}`$ is the co–moving, density–weighted volume over which the object can be observed within the sample limits in flux and solid angle. This luminosity function will reproduce the source counts of the input sample exactly.
We write the luminosity function as
$$\mathrm{\Phi }(L_\mathrm{x},z)=\mathrm{\Phi }(L_\mathrm{x},0)\rho (z,L_\mathrm{x})$$
(13)
where $`\rho (z,L_\mathrm{x})`$ is the space density, or density evolution, normalized to z=0. We approximate $`\rho (z,L_\mathrm{x})`$ by an analytical density template. For the redshift dependence we use at low redshifts a power law of $`(1+z)`$; at higher redshifts we adopt the shape of the density function used by Schmidt et al. (ssg95 (1995)) for optically selected quasars:
$`\rho _{\mathrm{tem}}(z)=(1+z)^m\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}0}<z<z_c`$ (14)
$`\rho _{\mathrm{tem}}(z)=(1+z_c)^mz_c<z<z_d`$ (15)
$`\rho _{\mathrm{tem}}(z)=(1+z_c)^m10^{k(zz_d)}z>z_d`$ (16)
We use the RBS which contains 205 AGN–1 in the $`\mathrm{log}L_\mathrm{x}`$ range 42–46, mostly at low redshifts, to derive the zero redshift luminosity function $`\mathrm{\Phi }(L_\mathrm{x},0)`$. The main excercise then is to derive the values of the template parameters $`(m,z_c,z_d,k)`$ by fitting to the flux and redshift distributions in the deeper X–ray surveys. This will allow the direct derivation of the luminosity function at the center of bins of luminosity and redshift, as outlined in the steps below.
1. Start the iteration by assuming initial values of the template parameters $`(m,z_c,z_d,k)`$ as a function of $`L_\mathrm{x}`$;
2. The zero redshift luminosity function is the sum of delta functions for the assembly of RBS sources
$$\mathrm{\Phi }(L_\mathrm{x},0)=\underset{i=1}{\overset{205}{}}(1/V_{\mathrm{max},i}^{\mathrm{RBS}})\delta (L_\mathrm{i}L_\mathrm{x})$$
(17)
where $`V_{\mathrm{max},i}^{\mathrm{RBS}}`$ is the accessible density–weighted volume of RBS source $`i`$ in the RBS, based on its solid angle and flux limit distribution;
3. Next, predict expected numbers and redshifts in the deeper surveys. For the part of the luminosity function based on RBS source $`i`$, the expected number in survey $`sur`$ based on the density template is
$$n_{\mathrm{tem}}(L_\mathrm{i},z)\mathrm{\Delta }z=(V_{\mathrm{max},i}^{sur}(z)/V_{\mathrm{max},i}^{\mathrm{RBS}})\mathrm{\Delta }z$$
(18)
where $`V_{\mathrm{max},i}^{sur}(z)`$ is the accessible density–weighted volume of RBS source $`i`$ over the redshift range $`z,z+\mathrm{\Delta }z`$ based on the solid angle and flux limit distributions in survey $`sur`$;
4. In order to compare with the observed numbers, we use four luminosity classes with $`\mathrm{\Delta }\mathrm{log}L_\mathrm{x}=1.0`$ centered on $`\mathrm{log}L_\mathrm{x}=42.5,43.5,44.5,45.5`$ and redshift shells of $`\mathrm{\Delta }\mathrm{log}z=0.1`$ centered on $`\mathrm{log}z=0.95,0.85,\mathrm{},0.75`$. Total binned numbers predicted in all the non–RBS surveys are $`N_{\mathrm{tem}}(L_\mathrm{x},z)`$;
5. Next, use the observed number $`N_{\mathrm{obs}}(L_\mathrm{x},z)`$ in each $`(L_\mathrm{x},z)`$ bin to derive the space density or luminosity function at the center of each luminosity bin by scaling the template value of the luminosity function:
$$\mathrm{\Phi }_{\mathrm{obs}}(L_\mathrm{x},z)=\mathrm{\Phi }(L_\mathrm{x},0)\rho _{\mathrm{tem}}(z,L_\mathrm{x})N_{\mathrm{obs}}(L_\mathrm{x},z)/N_{\mathrm{tem}}(L_\mathrm{x},z)$$
(19)
6. Comparison of $`\mathrm{\Phi }_{\mathrm{obs}}(z,L_\mathrm{x})`$ and $`\mathrm{\Phi }(L_\mathrm{x},0)\rho _{\mathrm{tem}}(z,L_\mathrm{x})`$ serves as a guide for the next iteration of the template parameters $`(m,z_\mathrm{c},z_\mathrm{d},k)`$, starting at step 2 above.
### 5.2 Deriving $`V_{\mathrm{max}}`$ values
The values of $`V_{\mathrm{max},i}^{sur}`$ are based on the maximum accessible redshifts $`z_{\mathrm{max},i}^{sur}`$ of RBS source $`i`$ in survey $`sur`$. For the SA–N, NEPS, RIXOS, RMS and CDF–S surveys, where missing redshifts are not correlated with X-ray or optical flux, we assume that $`z_{\mathrm{max},i}^{sur}`$ equals the X-ray redshift limit $`z_{\mathrm{lim},\mathrm{x}}`$. In accounting for the missing redshifts, we assume that their distribution is similar to that of the observed redshifts and use an effective solid angle appropriately reduced from the geometric solid angle.
For the RDS/XMM and CDF–N surveys, where redshifts are missing primarily for optically very faint objects, we proceed as follows. In this case, we employ only objects brighter than a magnitude limit $`m_{\mathrm{lim}}`$ (in practice $`R<24`$) and use the full geometric survey solid angle. In predicting the number of expected objects in survey $`sur`$, this introduces a second limiting redshift, $`z_{\mathrm{lim},\mathrm{opt},i}`$ that will depend on the optical luminosity of RBS source $`i`$ contributing to the luminosity function. The relevant limit $`z_{\mathrm{max},i}^{sur}`$ for the derivation of the expected counts is the smaller of $`z_{\mathrm{lim},\mathrm{x},i}`$ and $`z_{\mathrm{lim},\mathrm{opt},i}`$.
The determination of $`z_{\mathrm{lim},\mathrm{opt},i}`$ is based on the magnitudes of the RBS objects. Among the 205 RBS AGN–1, Salvato (salvato (2002)) has carried out optical photometry of a redshift-limited sample of 89 sources, deriving magnitudes for the nucleus, disk and spheroid. Comparison of the total magnitudes given by Salvato with the RBS magnitudes shows $`<R_{\mathrm{RBS}}R_{\mathrm{Sal}}>=0.5`$ with a dispersion of 0.9 mag. We apply this correction to the 116 remaining RBS sources. Based on the systematics of the Salvato magnitudes, we find that on the average, the absolute R–magnitude of the galaxy light is –21.3, that the disk light accounts for 56% of the galaxy light, and that for low–luminosity objects the nucleus contributes no less that 25% of the total light. This allows the, admittedly uncertain, derivation of the nuclear, disk and spheroid magnitudes of the remaining 116 RBS objects. We use the following values for the spectral index $`\alpha =d\mathrm{log}S/d\mathrm{log}\nu `$ of the nuclear, disk and spheroid components, estimated from spectral energy distributions given by Kinney (kinney (1996)): $`\alpha _{\mathrm{nuc}}=0.5`$, $`\alpha _{\mathrm{dis}}=1.0`$, $`\alpha _{\mathrm{sph}}=3.0`$. These values allow evaluation of the K-correction required to derive the redshift $`z_{\mathrm{lim},\mathrm{opt},i}`$ at which each of the 205 RBS sources reaches the spectroscopic magnitude limit $`R_{\mathrm{lim}}=24`$. We illustrate in Figure 6 the effect of the optical limit on the redshift limits for the CDF–N: for objects below the diagonal line, the redshift limit is that set by the optical limit.
### 5.3 The zero redshift X–ray luminosity function
The zero redshift luminosity function $`\mathrm{\Phi }(L_\mathrm{x},0)`$ is derived from the RBS sources as outlined in Sect 5.1, Eq. 17. Since the accessible volume $`V_{\mathrm{max},i}^{\mathrm{RBS}}`$ is density-weighted, this derivation requires information about the density template, discussed in the next section. At the low redshifts of the RBS objects (typically 0.1), this involves only the $`m`$ parameter of the low redshift $`(1+z)^m`$ density variation. A binned version of $`\mathrm{\Phi }(L_\mathrm{x},0)`$, created by summing the delta functions in bins of $`\mathrm{\Delta }\mathrm{log}L_\mathrm{x}=0.2`$, is shown in Fig. 7.
### 5.4 The analytical density template $`\rho (z,L_\mathrm{x})`$
The density template defined in Eqs. 1516 with parameters $`(m,z_\mathrm{c},z_\mathrm{d},k)`$ is an analytical approximation of the density function $`\rho (z,L_\mathrm{x})`$. The derivation of the template parameters was carried out as follows. As described in Sect. 5.1, it is a procedure of trial and error. For each of the four luminosity classes, in iterating the value of $`m`$, we minimize $`<N_{\mathrm{obs}}(L_\mathrm{x},z)N_{\mathrm{tem}}(L_\mathrm{x},z)>`$ for $`0.5<z<z_\mathrm{c}`$, since at low redshifts the luminosity function is firmly anchored by the RBS objects. The values of $`z_\mathrm{c}`$, $`z_\mathrm{d}`$ and $`k`$ are derived by minimizing $`<N_{\mathrm{obs}}(L_\mathrm{x},z)N_{\mathrm{tem}}(L_\mathrm{x},z)>`$ for $`z>z_\mathrm{c}`$, by adjusting $`k`$ for the lower luminosity classes to fit the total observed AGN–1 in the CDF–S and CDF–N (regardless of availability of redshifts or identification), and for the higher luminosity classes by adopting parameter values for optically selected quasars (see next section). This procedure makes it difficult to evaluate the statistical significance of the templates. We will show below that the accuracy of the templates actually has a negligible effect on the values of $`\mathrm{\Phi }_{\mathrm{obs}}(L_\mathrm{x},z)`$.
The resulting values of $`(m,z_\mathrm{c},z_\mathrm{d},k)`$ for the four luminosity classes are given in Table 6. The errors on for the low–redshift power index $`m=d\mathrm{log}\rho /d\mathrm{log}(1+z)`$ have been estimated using the number of observed objects in the redshift shells $`0.5<z<z_\mathrm{c}`$ in comparison with the RBS objects. From this analysis and the results shown in Table 4 it appears that $`m`$ exhibits a significant increase with $`L_\mathrm{x}`$ for $`\mathrm{log}L_\mathrm{x}>43.0`$. For these luminosities, we employ a quadratic interpolation of $`m`$ with $`\mathrm{log}L_\mathrm{x}`$ through the values given in the table. For $`\mathrm{log}L_\mathrm{x}=4243`$ we assume $`m`$ to be constant.
### 5.5 The X–ray luminosity function $`\mathrm{\Phi }_{\mathrm{obs}}(L_\mathrm{x},z)`$
The observed values of the luminosity function $`\mathrm{\Phi }_{\mathrm{obs}}(L_\mathrm{x},z)`$ are obtained by scaling the template luminosity function by a factor $`N_{\mathrm{obs}}(L_\mathrm{x},z)/N_{\mathrm{tem}}(L_\mathrm{x},z)`$, see Eq. 19. These values of $`\mathrm{\Phi }_{\mathrm{obs}}(L_\mathrm{x},z)`$ are plotted versus the redshift for each of the four luminosity classes in Figs. 89. The $`\pm 1\sigma `$ error bars are based on the numbers $`N_{\mathrm{tem}}(L_\mathrm{x},z)`$ predicted by the template. The figures also show the template luminosity function resulting from the product of the zero redshift luminosity function $`\mathrm{\Phi }(L_\mathrm{x},0)`$ and the template density function $`\rho _{\mathrm{tem}}(z,L_\mathrm{x})`$.
The $`\pm 1\sigma `$ error bars do not include any contribution reflecting the error of the template density function $`\rho _{\mathrm{tem}}(z,L_\mathrm{x})`$. We explore the effect of the template by rederiving the predicted number $`N_{\mathrm{tem}}(L_\mathrm{x},z)`$ on the extreme assumption that $`m=0`$ and $`k=0`$, i.e that there is zero density evolution. We find that the observed densities $`\mathrm{log}\mathrm{\Phi }_{\mathrm{obs}}(L_\mathrm{x},z)`$ increase by only $`0.000.06`$. Since the template errors must be much smaller than assumed in this extreme example, their effect on the error of $`\mathrm{\Phi }_{\mathrm{obs}}(L_\mathrm{x},z)`$ is negligible. The main result of this section is the set of $`\mathrm{\Phi }_{\mathrm{obs}}(L_\mathrm{x},z)`$ values plotted in Figs. 89; the templates serve primarily to eliminate the uncertainty related to binning.
For $`\mathrm{log}L_\mathrm{x}=4243`$ and $`4344`$, the density rises by an order of magnitude to $`z0.7`$ and $`1.2`$, respectively, and then declines steadily (Fig. 8). As shown in the next section, the peak at $`z0.7`$ for $`\mathrm{log}L_\mathrm{x}=4243`$ is little affected by redshift spikes caused by large-scale structure.
The density distribution for AGN–1 with $`\mathrm{log}L_\mathrm{x}`$=44–45 is documented to $`z4`$ (Fig. 9). Since in this case, the evidence for a decline in density at high redshift was not initially clear, we adopted for the parameters $`(z_\mathrm{c},z_\mathrm{d},k)`$ of the template above $`z=1.7`$ those found for high–luminosity optical quasars by Schmidt et al. (ssg95 (1995)). It appears from Fig. 9 that the X–ray data are consistent with the adopted shape. In order to further investigate whether the space density declines significantly at high redshift, we explore a test template in which the density does not decline at all, i.e. remains flat at high redshifts $`(k=0)`$. The results are shown in Fig. 10. The error bars are much reduced from those in Fig. 9, reflecting the fact that our errors are based on the predicted numbers, which are larger for a flat evolution function at high redshift. The bins at $`z=4.47,5.62,7.07`$ have observed/predicted numbers of 1/7.9, 0/7.2, and 0/19.3, respectively. Limiting the case to $`z<6.3`$, for which the observation of redshifts in high-luminosity quasars should be no problem, we have for $`z>4`$ one observed object versus 15 expected. The Poisson probability for such an occurence barring systematic effects is $`3\times 10^7`$, constituting strong evidence that for $`\mathrm{log}L_\mathrm{x}=4445`$ the space density declines beyond $`z4`$.
The low-redshift density parameter $`m`$ for $`\mathrm{log}L_\mathrm{x}=4546`$ continues the trend of an increasing $`m`$ for larger $`L_\mathrm{x}`$. At large redshift, the sparse data do not give any information about the density beyond $`z3`$ (see Fig. 9).
## 6 Discussion
In the present paper we have used two different methods to derive the AGN–1 X–ray luminosity function and its evolution. Detailed descriptions of the two methods are given in Sects. 4.1 and 5.1. Conceptually, the binned method derives a first order luminosity function by dividing the numbers $`N_{\mathrm{obs}}(L_\mathrm{x},z)`$ observed in the input surveys by the appropriate volumes. An analytical representation of the luminosity function is used to predict the numbers $`N_{\mathrm{mdl}}(L_\mathrm{x},z)`$ expected in these surveys. The luminosity function is then corrected by the factor $`N_{\mathrm{obs}}/N_{\mathrm{mdl}}`$.
In the $`V_{\mathrm{max}}`$ method, the RBS is used to derive the zero redshift luminosity function. An analytical density template is used to predict the numbers $`N_{\mathrm{tem}}(L_\mathrm{x},z)`$ expected in the deeper surveys. In this process, the effect of a spectroscopic magnitude limit on the $`V_{\mathrm{max}}`$ values is included for deep surveys where this limit applies. Once the template predictions are close to the observed numbers $`N_{\mathrm{obs}}(L_\mathrm{x},z)`$ in the deeper surveys, the luminosity function at the center of each $`(L_\mathrm{x},z)`$ bin is derived by multiplying the template luminosity function by $`N_{\mathrm{obs}}/N_{\mathrm{tem}}`$.
Using the optical magnitudes and spectra of the RBS sources in deriving the redshifts at which they would reach the spectroscopic limit $`R=24`$ introduces uncertainties. The magnitudes of the 116 RBS sources not studied by Salvato (salvato (2002)) are quite poor and so is their separation in nucleus, disk and spheroid components. The assumed spectra of these components are schematic. Obtaining spectral energy distributions for all the RBS sources down to the far UV would allow deriving their flux directly at or near the rest wavelengths corresponding to the $`R`$ magnitude at the limiting redshifts.
It is reassuring that the general properties and absolute values of the space density are very similar in the two different derivations in Sections 4 and 5. Figure 11 shows a direct comparison between the binned and $`V_{\mathrm{max}}`$ determinations of the space density, which agree very well within statistical errors.
We use overall fits to the luminosity function and its cosmological evolution in order to compare to previous work and to enable theoretical calculations with simplified analytical forms. For this purpose, we use two functional forms, i.e. a pure luminosity evolution (PLE) and a luminosity–dependent density evolution (LDDE) model. Assuming PLE, the luminosity evolution index $`p12.7`$ and the cutoff redshift $`z_\mathrm{c}1.7`$ are in rough agreement with previous results from type–1 AGNs in ROSAT or Einstein/ROSAT combined surveys (e.g. paper I; Jones et al.jones97 (1997); Page et al.page97 (1997); Boyle et al. boy93 (1993)).
The results of the K–S test for the PLE fit over the whole sample are marginally acceptable, with $`5\%`$ chance of obtaining a 2D K–S value larger than observed. However, this is caused by the sheer number of AGNs in the part of the z–L<sub>x</sub> space where PLE is still a good description, which dominates the overall statistics. Figure 5 clearly shows that PLE fails to reproduce the behavior of the SXLF around $`\mathrm{log}L_\mathrm{x}44`$ at $`0.4z1.7`$. The luminosity bin $`42<\mathrm{log}\mathrm{L}_\mathrm{x}<43.5`$ in the $`0.4<z<0.8`$ shell contains 41 objects where PLE predicts 19 and the adjacent regime of $`43.0<\mathrm{log}\mathrm{L}_\mathrm{x}<44.0`$ at $`0.8<z<1.2`$ contains 46 objects where PLE predicts 25, so that the overall excess corresponds to more than $`4.5\sigma `$. In constructing the LDDE form, we tried to fully represent the SXLF from our data with unprecedented redshift and luminosity coverage. The overall 2D K–S acceptance of LDDE has been improved to 36%. The only location where the LDDE model still deviates from the data significantly is the very end of our sample, the $`42<\mathrm{log}<43`$ bin at $`0.015<z<0.2`$, where LDDE predicts 66 objects while we observed 45 objects (a $`3\sigma `$ deviation).
Even though our sample is a soft X–ray-selected type–1 AGN sample, the overall behaviour of our XLF is similar to that obtained by Ueda et al. (ued03 (2003)) for the intrinsic (de–absorbed) luminosity function of hard X–ray selected obscured and unobscured AGN. To make this comparison, we have refitted our sample with the LDDE model where $`\beta _1`$ and $`\beta _2`$ are fixed to zero. This is exactly the same as the function form which Ueda et al. (ued03 (2003)) used to describe their intrinsic HXLF. All the $`z=0`$ XLF parameters and evolution parameters are remarkably close between our SXLF and HXLF, except the global normalization. The HXLF normalization is found to be about five times larger than that of our SXLF, after adjustments for the differences in energy bands and the difference in the luminosities at which normalizations are evaluated. This factor probably accounts for the absorbed objects missing in the SXLF. However, the Ueda et al. sample, containing about 250 AGN, is limited to lower luminosities and lower redshifts than our sample of 1000 objects, so that its statistical quality and limited sensitivity range were not sufficient to constrain the decline of the space density at high redshifts, which has been measured significantly in our sample for the first time.
Very recently, Barger et al., bar05 (2005) have presented X–ray luminosity function analyses both in the hard and soft X-ray bands, based on the CDF–N, CDF–S, CLASXS and ASCA surveys. Again, their results are in good agreement with the soft XLF discussed here and the hard XLF presented by Ueda et al., however, they still suffer from substantial identification incompleteness. Also, their results on broad-line AGN are not directly comparable to our type–1 AGN sample, because they only include the optically classified type–1 AGN and thus miss most of the low–luminosity unabsorbed AGN–1 we are concerned with in this paper. A critical comparison of our XLF with those of Ueda et al., ued03 (2003) and Barger et al., bar05 (2005) will be the topic of a future paper.
The evidence for a peak in the evolution function is quite strong at $`z0.7`$ for $`L_\mathrm{x}=4243`$ (see Figs. 5(a) and 8.
The faintest end of the sample depends on the small field of Chandra Deep Field-South and there is some concern on the effects of the large-scale structure and cosmic variance associated with it. In particular, Gilli et al. gil03 (2003) found two redshift spikes, one at $`z0.67`$ and the other at $`z0.73`$. One may wonder whether the z=0.7 peak is caused by these redshift spikes (see also Gilli et al., gil05 (2003)). These do, however, not affect our SXLF estimates significantly. In the analysis in Sect. 4 only six out of 41 AGNs (15%) in our sample in the $`0.4<z0.8`$ bin in the range $`42\mathrm{log}L_\mathrm{x}<43.5`$ are in these spikes ($`0.664z0.685`$ and $`0.725z0.742`$) and this is the only regime where the spikes give a non-negligible contribution. In the analysis in Sect. 5, the z=0.71 bin contains 5 objects from the two spikes, for a total of 5 out of 15 observed objects. Disregarding these objects would decrease the derived density by dex 0.22 in Fig. 8, actually leading to better agreement with the template densities. We therefore conclude that cosmic variance is not significantly affecting our results on the evolution of the space density.
We show in Figure 12 the AGN–1 space density as a function of cosmic time. We see dramatic changes with $`L_\mathrm{x}`$. For declining $`L_\mathrm{x}`$ as we move from high-luminosity AGN or quasars to Seyfert galaxies, the main formation of the objects occurred at later cosmic times. For $`\mathrm{log}L_\mathrm{x}>44`$ the density curve is similar to that for quasars. It is an intriguing question whether the observed dependence of $`m`$ on $`L_\mathrm{x}`$ is accompanied by a corresponding dependence on $`L_{\mathrm{opt}}`$. At $`\mathrm{log}L_\mathrm{x}<44`$ the AGN–1 are mostly Seyfert galaxies, for which there is no comparable optical evidence about their density curve.
These new results paint a dramatically different evolutionary picture for low–luminosity AGN compared to the high–luminosity QSOs. Obviously, the rare, high–luminosity objects can form and feed very efficiently rather early in the universe. Their space density declines by more than two orders of magnitude at redshifts below z=2. The bulk of the AGN, however, has to wait much longer to grow with and shows a decline of space density by less than a factor of 10 at redshifts below one. The late evolution of the low–luminosity Seyfert population is very similar to that which is required to fit the Mid–infrared source counts and background (Franceschini et al., fra02 (2002)) and also the bulk of the star formation in the Universe (Madau et al., mad98 (1998)), while the rapid evolution of powerful QSOs traces more the merging history of spheroid formation (Franceschini et al. fra99 (1999)).
This kind of anti–hierarchical Black Hole growth scenario is not predicted in most of the semi–analytic models based on Cold Dark Matter structure formation models (e.g. Kauffmann & Haehnelt kau00 (2000); Wyithe & Loeb wyi03 (2003)). This could indicate two modes of accretion and black hole growth with radically different accretion efficiency (see e.g. Duschl & Strittmatter dus02 (2002)). A self–consistent model of the black hole growth which can simultaneously explain the anti–hierarchical X-ray space density evolution and the local black hole mass function derived from the $`M_{\mathrm{BH}}\sigma `$ relation assuming two radically different modes of accretion has recently been presented by Merloni (mer04 (2004)).
Finally, we compare the space density of soft X–ray selected QSOs from our sample to the one of optically-selected QSOs at the most luminous end. The comparison is plotted in Fig. 13. The $`z<2`$ number density curve for optically selected QSOs ($`M_{b_\mathrm{J}}<26.0`$) is from the combination of the 2dF and 6dF QSO redshift surveys by Croom et al. (croom04 (2004)). The $`z>2.7`$ number densities from Schmidt, Schneider & Gunn (ssg95 (1995)) and Fan et al. (fan01 (2001)) have been originally given for $`h_{70}=5/7,\mathrm{\Omega }_\mathrm{m}=1,\mathrm{\Omega }_\mathrm{\Lambda }=0`$. Their data points have been converted to our default cosmology and the $`M_\mathrm{B}`$ threshold has been re-calculated with an assumed spectral index of $`\alpha _\mathrm{o}=0.79`$ ($`f_\nu \nu ^{\alpha _\mathrm{o}}`$), following e.g. Vignali et al. (vig03 (2003)). The plotted curve from Schmidt, Schneider & Gunn / Fan et al. is for $`M_\mathrm{B}<26.47`$ under these new assumptions. A small correction of densities due to the cosmology conversion causing redshift-dependent luminosity thresholds has also been made, assuming $`d\mathrm{\Phi }/d\mathrm{log}L_\mathrm{B}L_\mathrm{B}^{1.6}`$ (Fan et al. fan01 (2001)). The space density for the soft X–ray QSOs for the luminosity class $`44<\mathrm{log}L_\mathrm{x}<45`$ has been plotted both for the binned and $`V_{\mathrm{max}}`$ determination. The Croom et al., (croom04 (2004)) space density had to be scaled up by a factor of 16 in order to match the X-ray density at $`z2`$. The Schmidt, Schneider & Gunn / Fan et al. data points have been scaled by a factor of 40 to match the soft X–ray data at $`z=2.7`$ in the plot. There is relatively little difference in the density functions between the X–ray and optical QSO samples, although we have to keep in mind, that both the rise and the decline of the space density is varying with X–ray luminosity, so that this comparison can only be illustrative until larger samples of high–redshift X–ray selected QSOs will be available.
Very recently, Wall et al. wal05 (2005), have presented an update of the space density evolution of the Parkes quarter–Jansky sample of flat–spectrum radio sources. They basically confirm the rise and fall of the QSO population as now seen both in the optical and X-ray QSO populations.
## 7 Conclusions and outlook
We have merged the Chandra and XMM–Newton deep survey data with the whole body of previously identified ROSAT AGN samples. We have selected only the type–1 AGN in all samples and treated only the detections and X–ray fluxes in the 0.5–2 keV band. The different samples cover an unprecedented five orders of magnitude in flux limit and six orders of magnitude in survey solid angle between the ROSAT Bright and serendipitous surveys, the XMM–Newton Lockman Hole survey and the Chandra Deep Surveys. The sample comprises 944 identified AGN–1 objects and only 57 unidentified sources, which could be AGN–1, i.e. roundabout 1000 objects. The luminosity–redshift diagram is almost homegeneously filled with our sample objects. With this sample we arrive at the following conclusions:
1. The new Chandra and XMM–Newton sources are predominantly Seyfert galaxies at a median luminosity of $`10^{43}`$ erg s<sup>-1</sup> and a median redshift around 0.7 and push the determination of the X–ray luminosity and space density functions into so far unexplored parameter ranges of redshift and luminosity.
2. AGN–1 are by far the largest contributors to the soft X–ray selected samples. Their evolutional properties are responsible for the break in the total X–ray sources counts in the 0.5–2 keV band.
3. The soft X–ray luminosity function of AGN shows a clear change of shape as a function of redshift, confirming earlier reports of luminosity–dependent density evolution for optical quasars and X–ray AGN.
4. The space density function changes significantly for different luminosity classes. It shows a strong positive evolution, i.e. a density increase at low redshifts up to a certain redshift and then a flattening. The redshift, at which the evolution peaks, changes considerably with X–ray luminosity, from $`z0.50.7`$ for luminosities $`\mathrm{log}L_\mathrm{x}=4243`$ erg s<sup>-1</sup> to $`z2`$ for $`\mathrm{log}L_\mathrm{x}=4546`$ erg s<sup>-1</sup>.
5. The amount of density evolution from redshift zero to the maximum space density is also a strong function of X–ray luminosity. The change is more than a factor of 100 at high luminosities, similar to what has been observed for optically selected QSOs, but it is less than a factor of 10 for low X–ray luminosities.
6. For the first time, we find a clear decline of the space density of X–ray selected AGN towards high redshift, using a rigorous treatment of optical incompleteness and the corresponding survey volume. The decline is observed clearly for X–ray luminosities in the range $`\mathrm{log}L_{rmx}=4245`$ erg s<sup>-1</sup>, while at higher luminosities the survey volume at high–redshift is still too small to obtain meaningful densities.
In the future, X–ray surveys which are both wide and deep are necessary, in order to provide enough volume for a better measurement of the space density function of the rare high–luminosity AGN at large redshifts. Several new surveys towards this goal are already underway, e.g. the Chandra Multiwavelength Project (Champ) (Silverman et al., silver04 (2003)), the XMM-Newton AXIS project (Barcons et al., barc02 (2002)), the Chandra Large Area Synoptic X-Ray Survey (CLASXS) (Yang et al., yan04 (2004); Steffen et al., ste04 (2004)), the Extended Chandra Deep Field South (PI: W.N. Brandt) or the XMM-Newton COSMOS Field (PI: G. Hasinger), which together should enrich the sample of $`z>4`$ objects by about an order of magnitude. Ultimately, new X–ray Dark Energy missions, aiming to survey large solid angles on the sky to considerable depth could provide a factor of 100–1000 increase in AGN sample size.
###### Acknowledgements.
Part of this work was supported by the German *Deutsches Zentrum für Luft- und Raumfahrt, DLR* grant number 50 OR 0207. Also this work is partially supported by NASA grant NAG 5-10875 (LTSA). TM appreciates support from Max-Planck Society during his visits to MPE. We thank an anonymous referee for constructive comments, which significantly improved the paper.
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# 1 Introduction
## 1 Introduction
Let $`G`$ be an arbitrary graph on the vertex set $`V=\{v_1,\mathrm{},v_n\}`$. The edge cone of $`G`$ is the cone $`_+𝒜^n`$ spanned by the set $`𝒜`$ of all vectors $`e_i+e_j`$ such that $`v_i`$ is adjacent to $`v_j`$, where $`e_i`$ denotes the $`i`$th unit vector.
Our first main goal is to give an explicit combinatorial description of the edge cone of $`G`$, see Theorem 3.6 and Corollary 3.8. This description generalizes that of \[10, Corollary 3.3\]. In loc. cit. only the non bipartite case was studied.
The second main goal is to study in detail the facets of the edge cone of a connected bipartite graph and show a canonical irreducible representation of the edge cone, see Proposition 4.6 and Theorem 4.9. As an application the classical marriage theorem will follow. Our results can be applied to commutative algebra to compute the $`a`$-invariant and the canonical module of an edge subring along the lines of . It is a bit surprising that the edge cone of a bipartite graph has not been studied before from the point of view of polyhedral geometry.
To show our results we use graph theory, linear algebra (Farkas’s Lemma, incidence matrices of graphs, Carathéodory’s Theorem, Kronecker’s Lemma), and polyhedral geometry (finite basis theorem and facet structure). The proofs require a careful analysis at the graph theoretical level. Our main references for graphs, algebra and geometry are .
## 2 Preliminaries
Let $`𝒜=\{\alpha _1,\mathrm{},\alpha _q\}`$ be a finite set of vectors in $`^n`$. The cone $`_+𝒜`$ generated by the set $`𝒜`$ is defined as
$$_+𝒜:=\left\{\underset{i=1}{\overset{q}{}}a_i\alpha _i\right|a_i_+\text{ for all }i\}^n,$$
where $`_+`$ is the set of non negative real numbers.
By the finite basis theorem \[12, Chapter 4\] $`_+𝒜`$ is a rational polyhedral cone, that is, $`_+𝒜`$ is the intersection of finitely many closed halfspaces of the form:
$$H_a^+:=\{x^n|x,a0\},$$
where $`0a^n`$ and the nonzero entries of $`a`$ are relatively prime. Here $`x,a`$ denotes the standard inner product of $`x`$ and $`a`$.
Note that if $`H_a^{}:=H_a^+`$, then the intersection $`H_a^+H_a^{}`$ is the bounding hyperplane
$$H_a:=\{x|x,a=0\}$$
with normal vector $`a`$.
To simplify notation set $`Q=_+𝒜`$. Recall that a subset $`F^n`$ is a face of $`Q`$ if $`F=QH_a`$ for some hyperplane $`H_a`$ such that $`QH_a^+`$ or $`QH_a^{}`$. The hyperplane $`H_a`$ is called a supporting hyperplane of $`Q`$. The improper faces of $`Q`$ are $`Q`$ and $`\mathrm{}`$, all the other faces are called proper faces. If a face of $`Q`$ has dimension $`dim(Q)1`$ it is called a facet. The dimension of $`Q`$ is by definition the dimension of $`\mathrm{aff}(Q)`$, the affine hull of $`Q`$. Note that a face of $`Q`$ is again a finitely generated cone, see .
###### Definition 2.1
If a polyhedral cone $`Q=_+𝒜`$ in $`^n`$ is represented as
$$Q=\mathrm{aff}(Q)\left(\underset{i=1}{\overset{r}{}}H_{a_i}^+\right)$$
$`()`$
and satisfies
$$Q\mathrm{aff}(Q)\left(\underset{ij}{\overset{r}{}}H_{a_i}^+\right)$$
for all $`j`$, we say that $`()`$ is an irreducible representation of $`Q`$.
Part of the importance of an irreducible representation can be seen in the following general fact.
###### Theorem 2.2
Let $`Q`$ be a polyhedral cone in $`^n`$ which is not an affine space. If
$$Q=\mathrm{aff}(Q)H_{a_1}^+\mathrm{}H_{a_r}^+$$
is an irreducible representation of $`Q`$ with $`a_i^n\{0\}`$ for all $`i`$, then the facets of $`Q`$ are precisely the sets $`F_1,\mathrm{},F_r`$, where $`F_i=QH_{a_i}`$. Moreover each proper face of $`Q`$ is the intersection of those facets of $`Q`$ that contain it.
Proof. See \[12, Theorem 3.2.1\]. $`\mathrm{}`$
### The incidence matrix of a graph
In the sequel we use standard terminology and notation from graph theory and adopt the book of Harary as our main reference. For the reader’s convenience we recall a few notions about graphs.
Let $`G`$ be a simple graph with vertex set $`V(G)=\{v_1,\mathrm{},v_n\}`$ and edge set $`E(G)=\{z_1,\mathrm{},z_q\}`$, thus every edge $`z_i`$ is an unordered pair of distinct vertices $`z_i=\{v_{i_j},v_{i_k}\}`$. The incidence matrix $`M_G=[a_{ij}]`$ associated to $`G`$ is the $`n\times q`$ matrix defined by
$$a_{ij}=\{\begin{array}{cc}1\hfill & \text{if }v_iz_j,\text{ and}\hfill \\ 0\hfill & \text{if }v_iz_j.\hfill \end{array}$$
Note that each column of $`M_G`$ has exactly two $`1`$’s and the rest of its entries equal to zero. If $`z_i=\{v_{i_j},v_{i_k}\}`$ define $`\alpha _i=e_{i_j}+e_{i_k}`$, where $`e_i`$ is the $`i`$th unit vector in $`^n`$. Thus the columns of $`M_G`$ are precisely the vectors $`\alpha _1,\mathrm{},\alpha _q`$. As an example consider a triangle $`G`$ with vertices $`v_1,v_2,v_3`$. In this case:
$$M_G=\left(\begin{array}{ccc}1& 0& 1\\ 1& 1& 0\\ 0& 1& 1\end{array}\right),$$
with the vectors $`\alpha _1=e_1+e_2`$, $`\alpha _2=e_2+e_3`$, and $`\alpha _3=e_1+e_3`$ corresponding to the edges $`z_1=\{v_1,v_2\}`$, $`z_2=\{v_2,v_3\}`$, and $`z_3=\{v_1,v_3\}`$.
Recall that a graph $`G`$ is bipartite if there is a bipartition $`(V_1,V_2)`$ of $`G`$, that is, $`V_1`$ and $`V_2`$ are vertex classes satisfying:
(a) $`V(G)=V_1V_2`$,
(b) $`V_1V_2=\mathrm{}`$, and
(c) every edge of $`G`$ joins a vertex of $`V_1`$ to a vertex of $`V_2`$.
If $`G`$ is connected such a bipartition is uniquely determined. Equivalently $`G`$ is bipartite if all its cycles are of even length. If $`G`$ is bipartite, then its incidence matrix is totally unimodular, that is, all the $`i\times i`$ minors of $`M_G`$ are equal to $`0`$ or $`\pm 1`$ for all $`i1`$, see .
## 3 An explicit representation of the edge cone
Let us introduce some more terminology and fix some more notation that will be used throughout.
Let $`G`$ be a simple graph and let $`M_G`$ be its incidence matrix. We set $`𝒜_G`$ (or simply $`𝒜`$ if $`G`$ is understood) equal to the set $`\{\alpha _1,\mathrm{},\alpha _q\}`$ of column vectors of $`M_G`$. Since $`\alpha _i`$ represents an edge of $`G`$ sometimes $`\alpha _i`$ is called an edge or an edge vector. The edge cone of $`G`$ is defined as the cone $`_+𝒜`$ generated by $`𝒜`$. Note $`_+𝒜(0)`$ if $`G`$ is not a discrete graph. By one has
$$nc_0(G)=\mathrm{rank}(M_G)=dim_+𝒜,$$
where $`c_0(G)`$ is the number of bipartite connected components of $`G`$.
###### Lemma 3.1
If $`v_i`$ is not an isolated vertex of $`G`$, then the set $`F=H_{e_i}_+𝒜`$ is a proper face of the edge cone.
Proof. Note $`F\mathrm{}`$ because $`0F`$, and $`_+𝒜H_{e_i}^+`$. Since $`v_i`$ is not an isolated vertex $`_+𝒜H_{e_i}`$. $`\mathrm{}`$
Given a subset $`AV(G)`$, the neighbor set of $`A`$, denoted $`N_G(A)`$ or simply $`N(A)`$, is defined as
$$N(A)=\{vV(G)v\text{is adjacent to some vertex in}A\}.$$
Let $`A`$ be an independent set of vertices of $`G`$, that is, no two vertices of $`A`$ are adjacent. The supporting hyperplane of the edge cone of $`G`$ defined by
$$\underset{v_iA}{}x_i=\underset{v_iN(A)}{}x_i$$
will be denoted by $`H_A`$.
###### Lemma 3.2
If $`A`$ is an independent set of vertices of $`G`$ and $`F=_+𝒜H_A`$, then either $`F`$ is a proper face of the edge cone or $`F=_+𝒜`$.
Proof. It suffices to prove the containment $`_+𝒜H_A^{}`$. Take and edge $`\{v_j,v_{\mathrm{}}\}`$ of $`G`$. If $`\{v_j,v_{\mathrm{}}\}A\mathrm{}`$, then $`e_j+e_{\mathrm{}}`$ is in $`H_A`$, else $`e_j+e_{\mathrm{}}`$ is in $`H_A^{}`$. $`\mathrm{}`$
###### Definition 3.3
The support of a vector $`\beta =(\beta _i)^n`$ is defined as
$$\mathrm{supp}(\beta )=\{\beta _i|\beta _i0\}.$$
###### Lemma 3.4 ()
Let $`V=\{v_1,\mathrm{},v_n\}`$ be the vertex set of $`G`$ and let $`G_1,\mathrm{},G_r`$ be the connected components of $`G`$. If $`G_1`$ is a tree with at least two vertices and $`G_2,\mathrm{},G_r`$ are unicyclic non bipartite graphs, then $`\mathrm{ker}(M_G^t)=(\beta )`$ for some $`\beta `$ in $`^n`$ with $`\mathrm{supp}(\beta )=\{1,1\}`$ such that $`V(G_1)=\{v_iV|\beta _i=\pm 1\}`$.
For use below we recall the following form of Farkas’s Lemma, which is called the fundamental theorem of linear inequalities, see \[8, Theorem 7.1\].
###### Theorem 3.5
Let $`𝒜=\{\alpha _1,\mathrm{},\alpha _q\}`$ be a set of vectors in $`^n`$ and let $`\alpha ^n`$. If $`\alpha _+𝒜`$ and $`t=\mathrm{rank}\{\alpha _1,\mathrm{},\alpha _q,\alpha \}`$, then there exists a hyperplane $`H_a`$ containing $`t1`$ linearly independent vectors from $`𝒜`$ such that $`a,\alpha >0`$ and $`a,\alpha _i0`$ for $`i=1,\mathrm{},q`$.
###### Theorem 3.6
If $`G`$ is a connected graph with vertex set $`V=\{v_1,\mathrm{},v_n\}`$ and $`_+𝒜`$ is the edge cone of $`G`$, then
$$_+𝒜=\left(\underset{A}{}H_A^{}\right)\left(\underset{i=1}{\overset{n}{}}H_{e_i}^+\right),$$
$`()`$
where $``$ is the family of all the independent sets of vertices of $`G`$ and $`H_{e_i}^+`$ is the closed halfspace $`\{x^n|x_i0\}`$.
Proof. Let $`𝒜=\{\alpha _1,\mathrm{},\alpha _q\}`$ be the set of column vectors of the incidence matrix of $`G`$. Since $`_+𝒜`$ is clearly contained in the right hand side of Eq. $`()`$ it suffices to prove the other containment. Take $`\alpha ^n`$ in the right hand side of Eq. $`()`$. The proof is by contradiction, that is, assume $`\alpha _+𝒜`$. By \[10, Corollary 3.3\] we may assume $`G`$ bipartite with $`n3`$ vertices.
Note that if $`(V_1,V_2)`$ is the bipartition of $`G`$, then $`\mathrm{aff}(_+𝒜)`$ is the hyperplane
$$\underset{v_iV_1}{}x_i=\underset{v_iV_2}{}x_i,$$
because $`\mathrm{dim}(_+𝒜)=n1`$. As $`H_{V_1}^{}H_{V_2}^{}=H_{V_1}`$, the vector $`\alpha `$ is in $`\mathrm{aff}(_+𝒜)`$. As a consequence $`\mathrm{rank}(𝒜\{\alpha \})=n1`$.
By Theorem 3.5 there is $`a^n`$ and there are linearly independent vectors $`\alpha _1,\mathrm{},\alpha _{n2}`$ in $`𝒜`$ such that
(i) $`a,\alpha _i=0\text{ for }i=1,2,\mathrm{},n2`$,
(ii) $`a,\alpha _i0\text{ for }i=1,2,\mathrm{},q`$, and
(iii) $`a,\alpha >0`$.
Observe that $`_+𝒜H_a`$ because $`\mathrm{aff}(_+𝒜)H_a`$. There exists $`\alpha _j`$ in $`𝒜`$ such that $`\alpha _1,\mathrm{},\alpha _{n2},\alpha _j`$ is a basis of $`\mathrm{aff}(_+𝒜)`$ as a real vector space. In particular we can write
$$\alpha =\lambda _1\alpha _1+\mathrm{}+\lambda _{n2}\alpha _{n2}+\lambda _j\alpha _j(\lambda _i)$$
(1)
It follows that $`\alpha ,a=\lambda _j\alpha _j,a>0`$. Thus $`\lambda _j<0`$.
Consider the subgraph $`D`$ of $`G`$ whose edges correspond to $`\alpha _1,\mathrm{},\alpha _{n2}`$ and its vertex set is the union of the vertices in the edges of $`D`$. Set $`k=|V(D)|`$. By one has:
$$n2=\mathrm{rank}(M_D)=kc_0(D),$$
where $`M_D`$ is the incidence matrix of $`D`$ and $`c_0(D)`$ is the number of bipartite components of $`D`$. Thus $`0nk=2c_0(D)`$. This shows that either $`c_0(D)=1`$ and $`k=n1`$ or $`c_0(D)=2`$ and $`k=n`$.
Case (I): Assume $`C_0(D)=1`$ and $`k=n1`$. Set $`V(D)=\{v_1,\mathrm{},v_{n1}\}`$. As $`D`$ is a tree with $`n2`$ edges and $`\alpha _i,a=0`$ for $`i=1,\mathrm{},n2`$, applying Lemma 3.4, one may assume $`a=(a_1,\mathrm{},a_{n1},a_n)`$, where $`a_i=\pm 1`$ for $`1in1`$.
Set $`a^{}=(0,\mathrm{},0,1)=e_n`$. Next we prove the following
(a) $`\alpha _i,a^{}=0`$ for $`i=1,\mathrm{},n2`$.
(b) $`\alpha _j,a^{}=1`$ and $`\alpha _j,a<0`$.
Condition (a) is clear. To prove (b) first note $`\alpha _j(\alpha _1,\mathrm{},\alpha _{n2})`$. Then $`\alpha _j=e_k+e_n`$, because otherwise the “edge” $`\alpha _j`$ added to the tree $`D`$ form a graph with a unique even cycle, to derive a contradiction recall that a set of “edge vectors” forming an even cycle are linearly dependant. Thus $`\alpha _j,a^{}=1`$.
On the other hand $`\alpha _j,a<0`$, because if $`\alpha _j,a=0`$, then the hyperplane $`H_a`$ would contain the linearly independent vectors $`\alpha _1,\mathrm{},\alpha _{n2},\alpha _j`$ and consequently $`\mathrm{aff}(_+𝒜)`$ would be equal to $`H_a`$, a contradiction.
To finish the proof of this case we use the inequality
$$\alpha ,a^{}=\lambda _j\alpha _j,a^{}>0$$
to conclude $`\alpha ,a^{}>0`$, a contradiction because $`\alpha H_{e_n}^+`$.
Case (II): Assume $`C_0(D)=2`$ and $`k=n`$. Let $`D_1`$ and $`D_2`$ be the components of $`D`$ and set $`U_1=V(D_1)`$ and $`U_2=V(D_2)`$.
Using Lemma 3.4 we can relabel the vertices of the graph $`D`$ and write $`a=rb+sc`$, where $`0rs0`$ are rational numbers,
$$b=(b_1,\mathrm{},b_m,0,\mathrm{},0),c=(0,\mathrm{},0,c_{m+1},\mathrm{},c_n),$$
$`U_1=\{v_1,\mathrm{},v_m\}`$, $`b_i=\pm 1`$ for $`im`$, and $`c_i=\pm 1`$ for $`i>m`$. Set $`a^{}=b`$. Note the following:
(a) $`\alpha _i,a^{}=0`$ for $`i=1,\mathrm{},n2`$.
(b) $`\alpha _j,a^{}=1`$ and $`\alpha _j,a<0`$; this holds for any $`\alpha _j(\alpha _1,\mathrm{},\alpha _{n2})`$.
Condition (a) is clear. To prove (b) first note that the inequality $`\alpha _j,a<0`$ can be shown as in case (I). Observe that if an “edge” $`\alpha _k`$ has vertices in $`U_1`$ (resp. $`U_2`$), then $`\alpha _k,a=0`$. Indeed if we add the edge $`\alpha _k`$ to the tree $`D_1`$ (resp. $`D_2`$) we get a graph with a unique even cycle and this implies that $`\alpha _1,\mathrm{},\alpha _{n2},\alpha _k`$ are linearly dependant, that is, $`\alpha _k,a=0`$. Thus $`\alpha _j=e_i+e_{\mathrm{}}`$ for some $`v_iU_1`$ and $`v_{\mathrm{}}U_2`$. From the inequality
$$\alpha _j,a=r\alpha _j,b+s\alpha _j,c=rb_i+sc_{\mathrm{}}<0$$
we obtain $`b_i=1=\alpha _j,a^{}`$, as required.
Next we set
$$A=\{v_iV|b_i=1\}\text{ and }B=\{v_iV|b_i=1\}.$$
Note that $`\mathrm{}AU_1`$ and $`\mathrm{}BU_1`$, because $`D_1,D_2`$ are trees with at least two vertices. We will show that $`A`$ is an independent set of $`G`$ and $`B=N_G(A)`$.
If $`A`$ is not an independent set of $`G`$, there is an edge $`\{v_i,v_{\mathrm{}}\}`$ of $`G`$ for some $`v_i,v_{\mathrm{}}`$ in $`A`$. Thus $`\alpha _k=e_i+e_{\mathrm{}}`$, by (a) and (b) we get $`a^{},\alpha _k0`$, which is impossible because $`a^{},\alpha _k=2`$. This proves that $`A`$ is an independent set of $`G`$.
Next we show $`N_G(A)=B`$. If $`v_iN_G(A)`$, then $`\alpha _k=e_i+e_{\mathrm{}}`$ for some $`v_{\mathrm{}}`$ in $`A`$, using (a) and (b) we obtain $`a^{},\alpha _k=b_i+10`$ and $`b_i=1`$, hence $`v_iB`$. Conversely if $`v_iB`$, since $`D_1`$ has no isolated vertices, there is $`1kn2`$ so that $`\alpha _k=e_i+e_{\mathrm{}}`$, for some $`\mathrm{}`$, by (b) we obtain $`a^{},\alpha _k=1+b_{\mathrm{}}=0`$, which shows that $`v_{\mathrm{}}A`$ and $`v_iN_G(A)`$.
Therefore $`H_a^{}=H_A`$. Since $`_+𝒜H_a^{}`$ (this follows from (b)), there is $`\alpha _{\mathrm{}}H_A`$, thus $`H_a^{}^{}H_A^{}\mathrm{}`$ and consequently $`H_A^{}=H_a^{}^{}`$. By hypothesis $`\alpha H_A^{}`$, hence $`\alpha ,a^{}0`$. From Eq. (1) together with and (a) and (b) one has
$$\alpha ,a^{}=\lambda _j\alpha _j,a^{}=\lambda _j>0,$$
a contradiction. $`\mathrm{}`$
The next two results give an explicit representation by closed halfspaces of the edge cone of an arbitrary graph. Those representations were known for connected non bipartite graphs only .
###### Corollary 3.7
If $`G`$ is a graph with vertex set $`V=\{v_1,\mathrm{},v_n\}`$ and $`_+𝒜`$ is the edge cone of $`G`$, then
$$_+𝒜=\left(\underset{A}{}H_A^{}\right)\left(\underset{i=1}{\overset{n}{}}H_{e_i}^+\right),$$
$`()`$
where the intersection is taken over all the independent sets of vertices $`A`$ of $`G`$ and $`H_{e_i}^+=\{x^n|x_i0\}`$.
Proof. Let $`G_1,\mathrm{},G_r`$ be the connected components of $`G`$. For simplicity of notation we assume $`r=2`$ and $`V(G_1)=\{v_1,\mathrm{},v_m\}`$. There is a decomposition
$$_+𝒜=_+𝒜_{G_1}_+𝒜_{G_2}.$$
Let $`\delta `$ be a vector in the right hand side of Eq. $`()`$. One can write
$$\delta =(\delta _i)=\beta +\gamma =(\delta _1,\mathrm{},\delta _m,0,\mathrm{},0)+(0,\mathrm{},0,\delta _{m+1},\mathrm{},\delta _n).$$
Let $`A`$ be an independent set of $`G_1`$. Note $`N_G(A)=N_{G_1}(A)`$, hence
$$\underset{v_iA}{}\delta _i\underset{v_iN_G(A)}{}\delta _i=\underset{v_iN_{G_1}(A)}{}\delta _i.$$
Applying Theorem 3.6 yields $`\beta _+𝒜_{G_1}`$. Similarly one has $`\gamma _+𝒜_{G_2}`$. Hence $`\delta _+𝒜_G`$, as required. $`\mathrm{}`$
###### Corollary 3.8
Let $`G`$ be a graph with vertex set $`V=\{v_1,\mathrm{},v_n\}`$. Then a vector $`x=(x_1,\mathrm{},x_n)^n`$ is in $`_+𝒜`$ if and only if $`x`$ is a solution of the system of linear inequalities
$$\begin{array}{cccc}\hfill x_i& & 0,\hfill & i=1,\mathrm{},n\hfill \\ & & & \\ \hfill _{v_iA}x_i_{v_iN(A)}x_i& & 0,\hfill & \text{for all independent sets }AV.\hfill \end{array}$$
Proof. It follows at once from Corollary 3.7. $`\mathrm{}`$
###### Theorem 3.9
If $`G`$ is a graph with vertex set $`V=\{v_1,\mathrm{},v_n\}`$ and $`F`$ is a facet of the edge cone of $`G`$, then either
(a) $`F=_+𝒜\{x^n|x_i=0\}`$ for some $`1in`$, or
(b) $`F=_+𝒜H_A`$ for some independent set $`A`$ of $`G`$.
Proof. By Corollary 3.7 we can write
$$_+𝒜=\mathrm{aff}(_+𝒜)H_1^{}\mathrm{}H_r^{}$$
for some hyperplanes $`H_1,\mathrm{},H_r`$ such that none of the halfspaces $`H_j^{}`$ can be omitted in the intersection and each $`H_j`$ is either of the form $`H_{e_i}`$ or $`H_j=H_A`$ for some independent set $`A`$. By Theorem 2.2 the facets of $`_+𝒜`$ are precisely the sets $`F_1,\mathrm{},F_r`$, where $`F_i=H_i_+𝒜`$. $`\mathrm{}`$
###### Remark 3.10
(i) To verify whether a face $`F`$ as in (a) or (b) is a facet consider the set $``$ of all $`v_i𝒜`$ that are in $`F`$. Note that $`F`$ is a facet if and only if $`dim_+=r1`$, where $`r`$ is the dimension of the edge cone.
(ii) If $`G`$ is bipartite and connected, a graph theoretical characterization of the facets of the edge cone of $`G`$ will be given in Proposition 4.6. The facets of the edge cone of $`G`$ for $`G`$ non bipartite were characterized in \[10, Theorem 3.2\].
## 4 Studying the bipartite case
For connected bipartite graphs we will present sharper results on the irreducible representations of edge cones and give a characterization of their facets.
###### Proposition 4.1
Let $`G`$ be a connected bipartite graph with bipartition $`(V_1,V_2)`$. If $`A`$ is an independent set of $`G`$ such that $`AV_i`$ for $`i=1,2`$, then $`F=_+𝒜H_A`$ is a proper face of the edge cone.
Proof. Assume $`N(A)=V_2`$. Take any $`v_iV_1A`$ and any $`v_jV_2`$ adjacent to $`v_i`$, then $`e_i+e_jH_A`$. Thus we may assume $`N(A)V_i`$ for $`i=1,2`$.
Case (I): $`N(A)V_i\mathrm{}`$ for $`i=1,2`$. If the vertices in $`N(A)V_i`$ for $`i=1,2`$ are only adjacent to vertices in $`A`$, then pick vertices $`v_iN(A)V_i`$ and note that there is no path between $`v_1`$ and $`v_2`$, a contradiction. Thus there must be a vector in the edge cone which is not in $`H_A`$.
Case (II): $`AV_1`$. If the vertices in $`N(A)`$ are only adjacent to vertices in $`A`$. Then a vertex in $`A`$ cannot be joined by a path to a vertex in $`V_2N(A)`$, a contradiction. As before we obtain $`_+𝒜H_A`$. $`\mathrm{}`$
###### Proposition 4.2
Let $`G`$ be a connected bipartite graph with bipartition $`(V_1,V_2)`$ and $``$ the family of independent sets $`A`$ of $`G`$ such that $`H_A_+𝒜_G`$ is a facet. If $`A`$ is in $``$ and $`V_iA\mathrm{}`$ for $`i=1,2`$, then the halfspace $`H_A^{}`$ is redundant in the following expression of the edge cone
$$_+𝒜=\mathrm{aff}(_+𝒜)\left(\underset{A}{}H_A^{}\right)\left(\underset{i=1}{\overset{n}{}}H_{e_i}^+\right).$$
Proof. Set $`𝒜=\{\alpha _1,\mathrm{},\alpha _q\}`$. One can write $`A=A_1A_2`$ with $`A_iV_i`$ for $`i=1,2`$. There are $`\alpha _1,\mathrm{},\alpha _{n2}`$ linearly independent vectors in $`H_A_+𝒜`$, where $`n`$ is the number of vertices of $`G`$. Consider the subgraph $`D`$ of $`G`$ whose edges correspond to $`\alpha _1,\mathrm{},\alpha _{n2}`$ and its vertex set is the union of the vertices in those edges. Note that $`D`$ cannot be connected. Indeed there is no edge of $`D`$ connecting a vertex in $`N_G(A_1)`$ with a vertex in $`N_G(A_2)`$ because all the vectors $`\alpha _1,\mathrm{},\alpha _{n2}`$ satisfy the equation
$$\underset{v_iA}{}x_i=\underset{v_iN_G(A)}{}x_i.$$
Hence by the proof of Theorem 3.6 it follows that $`D`$ is a spanning subgraph of $`G`$ with two connected components $`D_1`$ and $`D_2`$ (which are trees) such that $`V(D_i)=A_iN_G(A_i)`$, $`i=1,2`$. Therefore $`H_{A_i}`$ is a proper support hyperplane defining a facet $`F_i=H_{A_i}_+𝒜`$, that is $`A_1,A_2`$ are in $``$. Since $`H_{A_1}^{}H_{A_2}^{}`$ is contained in $`H_A^{}`$ the proof is complete. $`\mathrm{}`$
###### Proposition 4.3
Let $`G`$ be a connected bipartite graph with bipartition $`(V_1,V_2)`$. If $`A_2V_2`$ and $`F=H_{A_2}_+𝒜`$ is a facet of the edge cone of $`G`$, then
$$H_{A_2}^{}\mathrm{aff}(𝒜^{})=\{\begin{array}{cc}H_{A_1}^{}\mathrm{aff}(𝒜^{})\hfill & \text{ where }A_1=V_1N(A_2)\mathrm{},\text{or}\hfill \\ H_{e_i}^+\mathrm{aff}(𝒜^{})\hfill & \text{for some vertex }v_i\text{ with }G\{v_i\}\text{ connected},\hfill \end{array}$$
where $`𝒜^{}=𝒜\{0\}`$.
Proof. Let us assume $`G`$ has $`p`$ vertices $`v_1,\mathrm{},v_p`$ and $`V_1`$ is the set of the first $`m`$ vertices of $`G`$. Set $`𝒜=\{\alpha _1,\mathrm{},\alpha _q\}`$. There are $`\alpha _1,\mathrm{},\alpha _{p2}`$ linearly independent vectors in the hyperplane $`H_{A_2}`$. Consider the subgraph $`D`$ of $`G`$ whose edges correspond to $`\alpha _1,\mathrm{},\alpha _{p2}`$ and its vertex set is the union of the vertices in those edges. As $`G`$ is connected either $`D`$ is a tree with $`p1`$ vertices or $`D`$ is a spanning subgraph of $`G`$ with two connected components.
If $`D`$ is a tree, write $`V(D)=V(G)\{v_i\}`$ for some $`i`$. Note
$$\alpha _j,\alpha _{A_2}=\alpha _j,e_i(j=1,\mathrm{},q),$$
where
$$\alpha _{A_2}=\underset{v_iA_2}{}e_i\underset{v_iN(A_2)}{}e_i.$$
Indeed if the “edge” $`\alpha _j`$ has vertices in $`V(D)`$, then both sides of the equality are zero, otherwise write $`\alpha _j=e_i+e_{\mathrm{}}`$. Observe $`v_iA_2`$ and $`v_{\mathrm{}}N(A_2)`$ because $`H_{A_2}`$ being a facet cannot contain $`\alpha _j`$, thus both sides of the equality are equal to $`1`$. As a consequence since $`\mathrm{aff}(𝒜^{})=(\alpha _1,\mathrm{},\alpha _{p2},\alpha _j)`$ for some $`\alpha _j=e_i+e_{\mathrm{}}`$ we rapidly obtain
$$\alpha ,\alpha _{A_2}=\alpha ,e_i(\alpha \mathrm{aff}(𝒜^{})).$$
Therefore
$$H_{A_2}^{}\mathrm{aff}(𝒜^{})=H_{e_i}^+\mathrm{aff}(𝒜^{}),$$
as required.
We may now assume $`D`$ is not a tree. We claim $`A_1=V_1N(A_2)\mathrm{}`$. If $`V_1=N(A_2)`$. Take $`v_iV_2A_2`$ and $`\{v_i,v_j\}`$ and edge of $`D`$ containing $`v_i`$. Hence since $`v_jN(A_2)`$ we get $`e_i+e_j,\alpha _{A_2}=1`$, a contradiction because $`e_i+e_j`$ is in $`H_{A_2}`$. Thus $`A_1\mathrm{}`$. Since all the vectors in $`\mathrm{aff}(𝒜^{})`$ satisfy the linear equation
$$\underset{i=1}{\overset{m}{}}x_i=\underset{i=m+1}{\overset{p}{}}x_i,$$
we obtain
$$H_{A_2}^{}\mathrm{aff}(𝒜^{})=\{x\mathrm{aff}(𝒜^{})|\underset{v_iV_1N(A_2)}{}x_i\underset{v_iV_2A_2}{}x_i\}.$$
Hence we need only show $`V_2A_2=N(A_1)`$. The containment $`N(A_1)V_2A_2`$ holds in general. For the reverse containment take $`v_iV_2A_2`$. There is $`v_j`$ such that $`\{v_i,v_j\}`$ is an edge of $`D`$. If $`v_jN(A_2)`$, then $`e_i+e_j,\alpha _{A_2}=1`$, a contradiction because $`e_i+e_jH_{A_2}`$. Hence $`v_jN(A_2)`$ and $`v_iN(A_1)`$. $`\mathrm{}`$
For later use we state the following duality of facets which follows from the proof of Proposition 4.3.
###### Lemma 4.4
Let $`G`$ be a connected bipartite graph with bipartition $`(V_1,V_2)`$ and let $`F=H_A_+𝒜`$ be a facet of $`_+𝒜`$ with $`AV_1`$. Then
(a) If $`N(A)=V_2`$, then $`A=V_1\{v_i\}`$ for some $`v_iV_1`$ and $`F=H_{e_i}_+𝒜`$.
(b) If $`N(A)V_2`$, then $`F=H_{V_2N(A)}_+𝒜`$ and $`N(V_2N(A))=V_1A`$.
###### Definition 4.5
For any set of vertices $`S`$ of a graph $`G`$, the induced subgraph $`S`$ is the maximal subgraph of $`G`$ with vertex set $`S`$.
###### Proposition 4.6
Let $`G`$ be a connected bipartite graph with bipartition $`(V_1,V_2)`$ and let $`AV_1`$. Then $`F=H_A_+𝒜`$ is a facet of $`_+𝒜`$ if and only if
(a) $`AN(A)`$ is connected with vertex set $`V(G)\{v\}`$ for some $`vV_1`$, or
(b) $`AN(A)`$ and $`(V_2N(A))(V_1A)`$ are connected and their union is a spanning subgraph of $`G`$.
Moreover any facet has the form $`F=H_A_+𝒜`$ for some $`AV_i`$, $`i=1`$ or $`i=2`$.
Proof. The first statement follows readily from Lemma 4.4 and using part of the proof of Theorem 3.6. The last statement follows combining Theorem 3.6 with Proposition 4.2. $`\mathrm{}`$
###### Remark 4.7
In Proposition 4.6 the case (a) is included in case (b). To see this make $`N(A)=V_2`$ and note that $`(V_2N(A))(V_1A)`$ must consist of a point. The condition in case (a) is equivalent to require $`G\{v\}`$ connected and in this case $`F=H_{e_i}_+𝒜`$, where $`v=v_i`$ correspond to the unit vector $`e_i`$.
###### Lemma 4.8
Let $`G`$ be a connected bipartite graph with bipartition $`(V_1,V_2)`$ and let $`F`$ be a facet of $`_+𝒜`$. If $`F=H_A_+𝒜=H_B_+𝒜`$ with $`AV_1`$ and $`BV_1`$, then $`A=B`$.
Proof. Set $`V_1=\{v_1,\mathrm{},v_m\}`$ and $`V_2=\{v_{m+1},\mathrm{},v_{m+n}\}`$. Recall that the equality
$$x_1+\mathrm{}+x_m=x_{m+1}+\mathrm{}+x_{m+n}$$
defines $`\mathrm{aff}(_+𝒜)`$.
Case (I): $`N(A)=V_2`$. Then by Lemma 4.4 (after permutation of vertices) $`A=\{v_1,\mathrm{},v_{m1}\}`$. Hence any $`xF`$ satisfies
$$\underset{v_iA}{}x_i\underset{v_iN(A)}{}x_i=x_m$$
and thus $`F=H_{e_m}_+𝒜`$. If $`v_mB`$, then $`\{v_m,v_j\}E(G)`$ for some $`v_j`$ in $`N(B)`$, thus $`e_m+e_jH_B`$ and consequently $`e_m+e_jH_{e_m}`$, a contradiction. Hence $`v_mB`$, that is, $`BA`$. If $`N(B)=V_2`$, then by Lemma 4.4 $`A=B`$. Assume $`V_2N(B)\mathrm{}`$, to complete the proof for this case we will show that this assumption leads to a contradiction. First note that $`v_m`$ is not adjacent to any $`v_jV_2N(B)`$. Indeed if $`\{v_m,v_j\}E(G)`$, then $`e_m+e_jH_B`$. Thus $`e_m+e_jH_{e_m}`$, a contradiction. Therefore by the connectivity of $`G`$ at least one vertex $`v_iV_1B`$ must be adjacent to both a vertex $`v_jV_2N(B)`$ and a vertex $`v_kN(B)`$, which is impossible because $`e_i+e_kH_{e_m}`$ and $`e_i+e_kH_B`$.
Case (II): $`N(A)V_2`$ and $`N(B)V_2`$. We begin by considering the subcase $`AB\mathrm{}`$. Take $`v_0AB`$ and $`v_0vB`$. By Proposition 4.6 the subgraph $`BN(B)`$ is connected, hence there is a path of even length
$$𝒫=\{v_0,v_1,v_2,\mathrm{},v_{2r1},v_{2r}=v\}$$
such that $`v_{2i}B`$ for all $`i`$. Note that $`v_2A`$. If $`v_2A`$, then $`e_1+e_2H_B`$ and $`e_1+e_2H_A`$, a contradiction. By induction we get $`v_{2i}AB`$ for all $`i`$. Hence $`vA`$. This proves $`BA`$, a similar argument proves $`A=B`$.
Assume now $`AB=\mathrm{}`$. We claim $`N(A)N(B)=\mathrm{}`$, for otherwise if $`\{v_j,v_k\}`$ is an edge with $`v_jB`$ and $`v_kN(A)N(B)`$, then $`e_j+e_kH_A`$ because $`v_jA`$ and $`e_j+e_kH_B`$, a contradiction.
We may now assume $`AB=N(A)N(B)=\mathrm{}`$. Observe $`ABV_2`$ because if $`V_2=AB`$, then $`G`$ would be disconnected with components $`AN(A)`$ and $`BN(B)`$. Take $`v_jV_1(AB)`$ such that $`v_j`$ is adjacent to some $`v_k`$ in $`N(A)N(B)`$, this choice is possible because $`G`$ is connected. Say $`v_kN(A)`$. Note $`v_kN(B)`$. Then $`e_j+e_kH_B`$ and $`e_j+e_kH_A`$, a contradiction. $`\mathrm{}`$
Putting together the previous results we obtain the following canonical way of representing the edge cone. The uniqueness follows from Lemma 4.4 and Lemma 4.8.
###### Theorem 4.9
If $`G`$ is a connected bipartite graph with bipartition $`(V_1,V_2)`$, then there is a unique irreducible representation
$$_+𝒜=\mathrm{aff}(_+𝒜)(_{i=1}^rH_{A_i}^{})(_iH_{e_i}^+)$$
such that $`A_iV_1`$ for all $`i`$ and $`v_iV_2`$ for $`i`$.
###### Lemma 4.10
If $`G`$ is a bipartite graph, then
$$^n_+𝒜=𝒜.$$
In particular if $`(\beta _1,\mathrm{},\beta _n)`$ is an integral vector in the edge cone, then $`_{i=1}^n\beta _i`$ is an even integer.
Proof. Let $`𝒜=\{\alpha _1,\mathrm{},\alpha _q\}`$ be the set of column vectors of the incidence matrix $`M`$ of $`G`$. Take $`\alpha ^n_+𝒜`$, then by Carathéodory’s Theorem \[4, Theorem 2.3\] and after an appropriate permutation of the $`\alpha _i`$’s we can write
$$\alpha =\eta _1\alpha _1+\mathrm{}+\eta _r\alpha _r(\eta _i0),$$
where $`r`$ is the rank of $`M`$ and $`\alpha _1,\mathrm{},\alpha _r`$ are linearly independent. On the other hand the submatrix $`M^{}=(\alpha _1\mathrm{}\alpha _r)`$ is totally unimodular because $`G`$ is bipartite (see ), hence by Kronecker’s lemma \[8, p. 51\] the system of equations $`M^{}x=\alpha `$ has an integral solution. Hence $`\alpha `$ is a linear combination of $`\alpha _1,\mathrm{},\alpha _r`$ with coefficients in $``$. It follows that $`\eta _i`$ for all $`i`$, that is, $`\alpha 𝒜`$. The other containment is clear. $`\mathrm{}`$
As an application we recover the following version of the marriage problem for bipartite graphs, see . Recall that a pairing off of all the vertices of a graph $`G`$ is called a perfect matching.
###### Theorem 4.11 (Marriage Theorem)
If $`G`$ is a bipartite graph, then $`G`$ has a perfect matching if and only if
$$|A||N(A)|$$
for every independent set of vertices $`A`$ of $`G`$.
Proof. Note that $`G`$ has a perfect matching if and only if the vector $`\beta =(1,1,\mathrm{},1)`$ is in $`𝒜`$. By Lemma 4.10 $`\beta `$ is in $`𝒜`$ if and only if $`\beta _+𝒜`$. Thus the result follows from Corollary 3.8. $`\mathrm{}`$
###### Corollary 4.12
Let $`G=𝒦_{m,n}`$ be the complete bipartite graph with $`mn`$. If $`V_1=\{v_1,\mathrm{},v_m\}`$ and $`V_2=VV_1`$ is the bipartition of $`G`$, then a vector $`z^{m+n}`$ is in $`_+𝒜`$ if and only if $`z=(x_1,\mathrm{},x_m,y_1,\mathrm{},y_n)`$ satisfies
$$\begin{array}{cccc}\hfill x_1+\mathrm{}+x_m& =& y_1+\mathrm{}+y_n,\hfill & \\ \hfill x_i& & 0,\hfill & i=1,\mathrm{},m,\hfill \\ \hfill y_i& & 0,\hfill & i=1,\mathrm{},n.\hfill \end{array}$$
In addition if $`m2`$, the inequalities define all the facets of $`_+𝒜`$.
###### Example 4.13
If $`G=𝒦_{1,3}`$ is the star with vertices $`\{v,v_1,v_2,v_3\}`$ and center $`x`$, then the edge cone of $`G`$ has three facets defined by
$$x_i0,(i=1,2,3).$$
Note that $`x=0`$, define a proper face of dimension $`1`$.
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# Single particle calculations for a Woods-Saxon potential with triaxial deformations, and large Cartesian oscillator basis
## 1 Purpose of the Fortran program
### 1.1 The Schrodinger equation
The program solves the Schrodinger equation for one body-deformed potential:
$$H|\varphi _i=E_i|\varphi _i$$
(1)
Here H represents the Hamiltonian of the system (neutrons or protons), and, $`E_i`$ , and, $`\varphi _i`$ , represent respectively, its eigenenergies, and its eigenfunctions.
### 1.2 The Hamiltonian
The Hamiltonian of the nucleon is defined by :
$$H=T+V+V^{SO}+e\varphi ^C$$
(2)
The quantities T, V, V<sup>so</sup>, indicate respectively, the kinetic, potential , and spin-orbit energy. For the proton, the Coulomb potential is represented by $`\varphi ^C`$, and $`e`$ is its charge
Explicitely:
$$T=\frac{\mathrm{}^2}{2m}.\stackrel{}{}^2$$
(3)
$`\mathrm{}`$ = Planck constant
m= nucleon masse
Owing to the fact that the nuclear forces have a short-range character, the average nuclear potential must ” follow on average” the nuclear density distribution:
For the case of the spherical symmetry (see also (19)) , we have:
$$V(r)\left(\frac{\rho __0}{1+\mathrm{exp}(\frac{rR}{a})}\right)=\frac{V__0}{1+\mathrm{exp}(\frac{rR}{a})}$$
(4)
For the deformed case, the above definition is generalized as :
$$V(\stackrel{}{r})=\frac{V_0}{1+\mathrm{exp}(R_VL_V/a_V)}$$
(5)
with
$`V_0,R_V,a_V`$ = mean field parameters
$`L_V`$ = quasi-radius (see eq. (9))
$$V^{so}(\stackrel{}{r})=\frac{1}{\mathrm{}}\left(\stackrel{}{}S(\stackrel{}{r})\stackrel{}{p}\right)\stackrel{}{\sigma }$$
(6)
with
$`\stackrel{}{p}=(\mathrm{}/i)\stackrel{}{}=`$ Neutron or proton momentum
$`\stackrel{}{\sigma }=(\sigma _x,\sigma _y,\sigma _z)`$ = Pauli spin-matrices
For the same reasons that for V, the mean field in the expression of the $`V^{so}`$ operator is given by a similar definition (see also the subsection 3.1) :
$$S(\stackrel{}{r})=\frac{\kappa }{1+\mathrm{exp}(R_{so}L_{so}/a_{so})}$$
(7)
with
$`\kappa `$ = spin-orbit coupling strength (there, the quantity $`S_0`$ is absorbed by $`\kappa `$, this latter is integrated to $`S(\stackrel{}{r})`$ and expressed in $`MeV.fm^2`$)
$`R_{so},a_{so}`$ = mean field parameters of $`S(\stackrel{}{r})`$
$`L_{so}`$ is the quasi-radius of the spin-orbit mean field (see eq. (12))
## 2 The Coulomb potential:
For the protons, the Coulomb’s potential is approximated by the one of the liquid drop model
$$\begin{array}{c}\mathrm{\Phi }^C(Z,P,\mathrm{\Phi })=\frac{\rho _{charge}}{4}\underset{Z_1}{\overset{Z_2}{}}dz\times \hfill \\ \hfill \underset{0}{\overset{2\pi }{}}𝑑\phi \left[\frac{(zZ)\frac{\rho _S^2}{z}+2\rho _S^22\rho _SP\mathrm{cos}(\phi \mathrm{\Phi })2P\frac{\rho _S}{\phi }\mathrm{sin}(\phi \mathrm{\Phi })}{\sqrt{(zZ)^2+\rho _S^2+P^22\rho _SP\mathrm{cos}(\phi \mathrm{\Phi })}}\right]\end{array}$$
(8)
where:
$`(Z,P,\mathrm{\Phi })`$ = cylindrical coordinates of the point where the Coulomb potential is calculated.(here, $`Z`$ must not be confused with the protons number)
$`\rho _{charge}`$ = $`(\stackrel{´}{Z}1)e/(4/3)\pi R_{ch}^3`$= charge density of the liquid drop
$`(\stackrel{´}{Z}1)`$= ”number of protons in the liquid drop”
$`R_{ch}`$= radius of the charge density
The integration domain is defined by the volume limited by the surface $`\pi _{ch}=0`$ (see eq. (15)).
The “nuclear surface of the protons” is given by $`\rho _S=\rho _{Surface}`$ in the eq. (16).
In fact, the code computes directly the quantity $`e\mathrm{\Phi }^C`$ (which is the Coulomb energy of the proton in the Coulomb field) instead the Coulomb potential $`\mathrm{\Phi }^C.`$( see the function *ephi*. in the code).
## 3 Supplementary details on the deformation of the mean field, and the different nuclear surfaces:
### 3.1 The quasi-radius and the nuclear surfaces
In the central average potential, the information on the distortion of the nuclear surface is given by the dimensionless quasi-radius $`L_V(\stackrel{}{r})`$. which is defined as :
$$L_V(\stackrel{}{r})=\frac{\mathrm{\Pi }_V(\stackrel{}{r})}{R_V\stackrel{}{}\mathrm{\Pi }_V(\stackrel{}{r})}$$
(9)
The quantity $`\mathrm{\Pi }_V(\stackrel{}{r})`$ is defined so that to recover the well-known spherical case (see section 4).
$$\mathrm{\Pi }_V(\stackrel{}{r})=\sqrt{\pi _V(\stackrel{}{r})\pi _{V\mathrm{min}}}\sqrt{\pi _{V\mathrm{min}}}$$
(10)
Here, $`\pi _{V\mathrm{min}}`$ is the absolute minimum of $`\pi _V(\stackrel{}{r})`$. This latter describes an hypersurface which is not the real nuclear surface, because generally $`\pi _V(\stackrel{}{r})0`$ in the expression of the quasi radius. The actual nuclear surface may be obtained by putting $`\pi _V(\stackrel{}{r})=0.`$
In this work, we have restricted ourselves only to simple ellipsoidal (triaxial) shapes for the effective nuclear surface.
$$\pi _V(\stackrel{}{r})=\frac{x^2}{A_V^2}+\frac{y^2}{B_V^2}+\frac{z^2}{C_V^2}1=0$$
(11)
$`A_V,B_V,`$ and, $`C_V`$ are thus the semi-axes of the ellipsoid.
In fact, we have to consider three distinct interactions, i.e. the central, the spin-orbit, and the Coulomb interaction. Therefore, we must define three respective surfaces (equations (11),(14),(15)). Thus, the equation (11) describes the nuclear surface relatively to the central interaction $`V(\stackrel{}{r})`$.
In completely analogous way, we have to define similar quantities for the spin-orbit interaction
$$L_{so}(\stackrel{}{r})=\frac{\mathrm{\Pi }_{so}(\stackrel{}{r})}{R_{so}\stackrel{}{}\mathrm{\Pi }_{so}(\stackrel{}{r})}$$
(12)
with,
$$\mathrm{\Pi }_{so}(\stackrel{}{r})=\sqrt{\pi _{so}(\stackrel{}{r})\pi _{so\mathrm{min}}}\sqrt{\pi _{so\mathrm{min}}}$$
(13)
Where, $`\pi _{so\mathrm{min}}`$ is the absolute minimum of $`\pi _{so}(\stackrel{}{r})`$, and the effective ”spin-orbit surface” is written as:
$$\pi _{so}(\stackrel{}{r})=\frac{x^2}{A_{so}^2}+\frac{y^2}{B_{so}^2}+\frac{z^2}{C_{so}^2}1=0$$
(14)
For the Coulomb potential, the effective nuclear surface is defined in the same way:
$$\pi _{ch}(\stackrel{}{r})=\frac{x^2}{A_{ch}^2}+\frac{y^2}{B_{ch}^2}+\frac{z^2}{C_{ch}^2}1=0$$
(15)
Following the expression of the equation (8), the Coulomb potential must be expressed in cylindrical coordinates Therefore, the equation of the effective ”Coulomb nuclear surface” (15) can be rewritten as:
$$\rho _{surface}^2=\frac{1\left(z/C_{ch}\right)^2}{\left\{\left(\mathrm{cos}\phi /A_{ch}\right)^2+\left(\mathrm{sin}\phi /B_{ch}\right)^2\right\}}$$
(16)
where:
$`x=\rho _{surface}\mathrm{cos}\phi `$ , $`y=\rho _{surface}\mathrm{sin}\phi `$ , $`z=z`$
The ”three densities”(neutrons, protons, spin-orbit) differ very little from each other. Therefore, the three surfaces are homothetic and slightly different to each other. Nevertheless, in order to simplify the problem, the protons distribution is assumed to be uniform in the calculation of the Coulomb potential.
### 3.2 The deformation parameters
Since we have three similar surfaces, and, so as to avoid repeating three times the same thing, we will omit to specify the indices of the surfaces. For example, the three volume conservation conditions are simply replaced by only one equation :
$$\frac{4}{3}\pi ABC=\frac{4}{3}\pi R^3$$
(17)
Actually, because of this condition, only two parameters are necessary
As usual, one will prefer the Bohr parameters $`(\beta ,\gamma )`$ instead of those of the ellipsoid.
$`A`$ $`={\displaystyle \frac{R}{\chi }}\left[1+\beta \left({\displaystyle \frac{5}{4\pi }}\right)^{1/2}\mathrm{cos}(\gamma {\displaystyle \frac{2}{3}}\pi )\right]`$ (18a)
$`B`$ $`={\displaystyle \frac{R}{\chi }}\left[1+\beta \left({\displaystyle \frac{5}{4\pi }}\right)^{1/2}\mathrm{cos}(\gamma {\displaystyle \frac{4}{3}}\pi )\right]`$ (18b)
$`C`$ $`={\displaystyle \frac{R}{\chi }}\left[1+\beta \left({\displaystyle \frac{5}{4\pi }}\right)^{1/2}\mathrm{cos}(\gamma )\right]`$ (18c)
$`R`$ is the radius, and, $`\chi `$ insures the volume conservation condition (17),
## 4 The Spherical case. The mean field parameters
### 4.1 The case of spherical symmetry
When A=B=C , or when the Bohr parameter $`\beta =0`$, the nuclear surface becomes spherical, and $`R_V`$, $`R_{so}`$, or, $`R_{ch}`$ represents simply the nucleus radius. In this case, we obtain the familiar Fermi-function for the two mean potentials ((5),(7)):
$$V(\stackrel{}{r})=\frac{V_0}{1+\mathrm{exp}\left[\left(rR_V\right)/a_V\right]}S(\stackrel{}{r})=\frac{\kappa }{1+\mathrm{exp}\left[\left(rR_{so}\right)/a_{so}\right]}$$
(19)
i.e. potentials of Woods-Saxon type.
The Coulomb’s potential (8) reduces to the well-known form:
| $`\mathrm{\Phi }^c(r)=\left[(Z1)e/2R_{ch}\right]\left[3(r/R_{ch})^2\right]`$ if$`rR_{ch}`$ |
| --- |
| $`\mathrm{\Phi }^c(r)=(Z1)e/r`$ if $`rR_{ch}`$ |
(20)
The spin-orbit interaction (6) can be expressed in the spherical case as:
$$V^{so}=(\frac{S(r)}{r}\frac{\stackrel{}{r}}{r}\stackrel{}{p})\stackrel{}{\sigma }=\frac{1}{r}\frac{S(r)}{r}(\stackrel{}{r}\stackrel{}{p})\stackrel{}{\sigma }$$
(21)
Finally, the $`V^{so}`$ operator takes the familiar form,
$$V^{so}(\stackrel{}{r})=\frac{1}{r}\frac{S(r)}{r}\stackrel{}{l}.\stackrel{}{\sigma }$$
(22)
The relations (19-22) of the spherical case are not used explicitly in the code. However, the well known spherical degeneracy of the energy levels were used in order to check the program. Moreover, the relation (20) serves as a first checking for the Coulomb potential.
### 4.2 Mean field parameters
Two options have been included in the code in order to choose between a particular set of parameters, or the Myers parameters . Thus, it is possible to define its own parameters in a separate file, or to employ those of Myers. In this latter case, the calculations are made in a suitable subroutine. In fact, only a part of the parameters set, namely, $`V_0,R_V,R_{so},R_{ch}`$ , is deduced from the droplet model of Myers, the remaining, i.e. $`\kappa ,a_{V\text{,}}a_{so}`$, are extracted from the Ref. . The explicit expressions for these parameters are given in appendix.
## 5 Principle of resolution
### 5.1 The method
The principle of this method consists to look for the eigenfunctions of the Schrodinger equation by their expansion on a truncated basis of the harmonic oscillator. In other words, the method used in solving such problem amounts essentially to writing the representative matrix of the Hamiltonian in this basis.
In practice, this method is characterized by two distinct stages. First, one builds the representative matrix of the Hamiltonian by means of the cited basis. Next, one diagonalizes this matrix in order to obtain the eigenvalues and the eigenvectors.
In our work, the cartesian coordinates are the most suitable.
### 5.2 The harmonic oscillator basis
The basis functions of the harmonic oscillator are defined as:
$$|n_xn_yn_z\mathrm{\Sigma }i^{ny}\varphi _{n_x}(x).\varphi _{n_y}(y).\varphi _{n_z}(z).\stackrel{}{\sigma }_\mathrm{\Sigma }$$
(23)
There, $`i^{n_y}`$ is a phase factor which insures, in accordance with the imposed symmetries (see section 6), that the matrix elements are real.
Explicitly:
$$\varphi _{n_x}(x)=\sqrt{\beta _x}\mathrm{exp}\left[\left(\beta _xx\right)^2/2\right].h_{n_x}(\beta _xx)with\beta _x=\sqrt{\frac{m\omega _x}{\mathrm{}}}$$
(24)
with analogous expressions for the y, and z axes.
The intrinsic spin states are:
$$\stackrel{}{\sigma }_{+1/2}=\left[\begin{array}{c}1\\ 0\end{array}\right],\stackrel{}{\sigma }_{1/2}=\left[\begin{array}{c}0\\ 1\end{array}\right]$$
(25)
The $`h_{n_x}`$, (or $`h_{n_y},`$ or $`h_{n_z}`$) quantities symbolize the normalized Hermite polynomials.
$$h_{n_x}(x)=H_{n_x}(x)/\sqrt{\left(2^{n_x}.n_x!\pi ^{1/2}\right)}$$
(26)
$`H_{n_x}(x)`$ are the usual Hermite polynomials
The quantum numbers $`n_x`$ , $`n_y`$ , $`n_z`$ , are integers, and give the order of the Hermite polynomials, $`\mathrm{\Sigma }=\pm \frac{1}{2}`$ represents the projection of the intrinsic spin on the z-axis.
At last, $`m`$ and $`\omega _x`$ , $`\omega _y`$ , $`\omega _z`$ represents the mass of the oscillator, i.e. the mass of the nucleon, and, its frequencies.
For convenient, the quantities $`\mathrm{}\omega _x`$, $`\mathrm{}\omega _y`$, $`\mathrm{}\omega _z`$,(or respectively $`\beta _x`$,$`\beta _y`$,$`\beta _z`$) are called the deformation parameters of the basis.
If the three frequencies are equal, the oscillator is then isotropic, and it can be characterized by only one frequency ($`\omega _x=\omega _y=\omega _z=\omega _0`$)
Furthermore, we assume that the nuclear surface of the oscillator is an equipotential. This involves the following condition:
$$\omega _x\omega _y\omega _z=\omega _0^3or\left(\mathrm{}\omega _x\right).\left(\mathrm{}\omega _y\right).\left(\mathrm{}\omega _z\right)=\left(\mathrm{}\omega _0\right)^3$$
(27)
### 5.3 The representative matrix of the Hamiltonian
With help of this basis, the matrix elements of the Hamiltonian $`H`$ can be written as:
$$n_x^{}n_y^{}n_z^{}\mathrm{\Sigma }^{}\left|H\right|n_xn_yn_z\mathrm{\Sigma }=n_x^{}n_y^{}n_z^{}\mathrm{\Sigma }^{}\left|T+V+V^{so}+e.\mathrm{\Phi }^c\right|n_xn_yn_z\mathrm{\Sigma }$$
(28)
#### 5.3.1 Matrix elements of the mean field V and the Coulomb energy e.$`\mathrm{\Phi }^C`$
Since $`V`$ does not depend on the spin, we adopt the following convenient notation for V:
$$n_x^{}n_y^{}n_z^{}\mathrm{\Sigma }^{}\left|V\right|n_xn_yn_z\mathrm{\Sigma }=\left(n_x^{}n_y^{}n_z^{}\left|V\right|n_xn_yn_z\right)\times \delta _{\mathrm{\Sigma }^{}\mathrm{\Sigma }}$$
(29)
where:
$$\begin{array}{c}\left(n_x^{}n_y^{}n_z^{}\left|V\right|n_xn_yn_z\right)\hfill \\ \hfill =i^{(n_yn_y^{})}\varphi _{n_x^{}}(x)\varphi _{n_y^{}}(y)\varphi _{n_z^{}}(z)V(x,y,z)\varphi _{n_x}(x)\varphi _{n_y}(y)\varphi _{n_z}(z)𝑑x𝑑y𝑑z=\\ \hfill =i^{(n_yn_y^{})}\underset{\mathrm{}}{\overset{+\mathrm{}}{}}\underset{\mathrm{}}{\overset{+\mathrm{}}{}}\underset{\mathrm{}}{\overset{+\mathrm{}}{}}e^{(x^2+y^2+z^2)}h_{n_x^{^{}}}(x).h_{n_y^{^{}}}(y).h_{n_z^{^{}}}(z)\\ \hfill \times V(\frac{x}{\beta _x},\frac{y}{\beta _y},\frac{z}{\beta _z}).h_{n_x}(x).h_{n_y}(y).h_{n_z}(z).dxdydz\end{array}$$
(30)
with $`\beta _x`$, $`\beta _y`$, $`\beta _z`$defined in eq.(24)
* Due to the parity of $`V`$ , the integral (30) vanishes if one of the three following conditions is not fulfilled:
$`(1)^{n_x}=(1)^{n_x^{}}`$
$`(1)^{n_y}=(1)^{n_y^{}}`$
$`(1)^{n_z}=(1)^{n_z^{}}`$
Therefore, it is not necessary to calculate them.
Since the respective indices must have the same parity, the complex factor of (30) can be rewritten as:
$`i^{(n_y^{}n_y)}=(1)^{(n_y^{}n_y)/2}`$
* Although the $`\left(n_x^{}n_y^{}n_z^{}\left|V\right|n_xn_yn_z\right)`$ elements are spin independent, they are stored actually as:$`n_x^{}n_y^{}n_z^{}\mathrm{\Sigma }^{}\left|V\right|n_xn_yn_z\mathrm{\Sigma }`$ in the computer memory
* The matrix elements of the Coulomb energy are calculated in the same way as those of $`V(\stackrel{}{r})`$ (putting $`e\varphi ^c(\stackrel{}{r})`$ instead $`V(\stackrel{}{r})`$), with the same change of scale.
* In the gaussian integration, we have to calculate the Hermite polynomials, the central mean potential $`V(\stackrel{}{r})`$, the spin-orbit mean potential $`S(\stackrel{}{r})`$, and the Coulomb potential $`\varphi ^c(\stackrel{}{r})`$ only at the nodes. Therefore, it is more convenient to store these quantities in specific arrays before any calculations (see the common/tabh/… declaration in the subroutine setsub).
#### 5.3.2 Matrix elements of the spin-orbit energy operator V<sup>so</sup>
Due to the presence of the derivative of $`S(\stackrel{}{r})`$, the direct calculation of the matrix elements of $`V^{so}`$ , i.e. $`n_x^{}n_y^{}n_z^{}\mathrm{\Sigma }^{}\left|V^{so}\right|n_xn_yn_z\mathrm{\Sigma }`$, is not convenient. These derivatives can be transferred on the basis functions by partial integration. Therefore, the derivatives of the basis functions can be expressed from the basis functions themselves by mean of the recursion relations. Finally, the matrix elements of $`V^{so}`$ can be obtained by suitable combinations of $`\left(n_x^{}n_y^{}n_z^{}\left|S\right|n_xn_yn_z\right)`$ elements, i.e.
$$\begin{array}{c}n_x^{}n_y^{}n_z^{}\mathrm{\Sigma }^{}\left|V^{so}\right|n_xn_yn_z\mathrm{\Sigma }=\hfill \\ \hfill \frac{m\omega _0}{2\mathrm{}}\left[2B_z\left(\mathrm{\Sigma }^{}\left|\sigma _z\right|\mathrm{\Sigma }\right)+B_+\left(\mathrm{\Sigma }^{}\left|\sigma _{}\right|\mathrm{\Sigma }\right)+B_{}\left(\mathrm{\Sigma }^{}\left|\sigma _+\right|\mathrm{\Sigma }\right)\right]\end{array}$$
(31)
where:
$$\sigma _\pm =\sigma _x\pm \sigma _y$$
(32)
$$B_\pm =B_xB_y$$
(33)
$$\begin{array}{c}B_x=\sqrt{\frac{\omega _y\omega _z}{\omega _0^2}}[\sqrt{n_y^{}(n_z+1)}(n_x^{}n_y^{}1,n_z^{}|S|n_xn_yn_z+1)\hfill \\ \hfill \sqrt{n_y(n_z^{}+1)}(n_x^{}n_y^{}n_z^{}+1\left|S\right|n_xn_y1,n_z)\\ \hfill +\sqrt{n_z(n_y^{}+1)}(n_x^{}n_y^{}+1,n_z^{}\left|S\right|n_xn_yn_z1)\\ \hfill +\sqrt{n_z^{}(n_y+1)}(n_x^{}n_y^{},n_z^{}1|S|n_xn_y+1,n_z)]\end{array}$$
(34)
$$\begin{array}{c}B_y=\sqrt{\frac{\omega _z\omega _x}{\omega _0^2}}[\sqrt{n_z^{}(n_x+1)}(n_x^{},n_y^{}n_z^{}1|S|n_x+1,n_yn_z)\hfill \\ \hfill +\sqrt{n_z(n_x^{}+1)}(n_x^{}+1,n_y^{}n_z^{}\left|S\right|n_xn_yn_z1)\\ \hfill \sqrt{n_x(n_z^{}+1)}(n_x^{}n_y^{}n_z^{}+1\left|S\right|n_x1,n_y,n_z)\\ \hfill +\sqrt{n_x^{}(n_z+1)}(n_x^{}1,n_y^{}n_z^{}|S|n_xn_yn_z+1)]\end{array}$$
(35)
$$\begin{array}{c}B_z=\sqrt{\frac{\omega _x\omega _y}{\omega _0^2}}[\sqrt{n_x^{}(n_y+1)}(n_x^{}1,n_y^{}n_z^{}|S|n_xn_y+1,n_z)\hfill \\ \hfill \sqrt{n_x(n_y^{}+1)}(n_x^{}n_y^{}+1n_z^{}\left|S\right|n_x1,n_yn_z)\\ \hfill +\sqrt{n_y(n_x^{}+1)}(n_x^{}+1,n_y^{}n_z^{}\left|S\right|n_xn_y1,n_x)\\ \hfill +\sqrt{n_y^{}(n_x+1)}(n_x^{}n_y^{}1,n_x^{}\left|S\right|n_x+1,n_yn_x)\end{array}$$
(36)
where $`S(\stackrel{}{r})`$ is given by eq.(7)
* The changes of sign in $`B_y`$ are involved by the phase factor of the basis functions.
* The computations of the matrix elements of $`S(\stackrel{}{r})`$ are carried out like those of $`V(\stackrel{}{r})`$) in (30).
#### 5.3.3 <br>Matrix elements of the kinetic energy operator T
At last, the matrix elements of the $`T`$ operator can be calculated in straightforward way:
$`n_x^{}n_y^{}n_z^{}\mathrm{\Sigma }^{}\left|T\right|n_xn_yn_z\mathrm{\Sigma }`$ $`={\displaystyle \frac{1}{4}}\delta _{\mathrm{\Sigma }\mathrm{\Sigma }^{}}i^{(n_yn_y^{})}\times `$
$`[\mathrm{}\omega _z\delta _{n_xn_x^{}}\delta _{n_yn_y^{}}\delta _{n_zn_z^{}}\text{ }(2n_z+1)`$
$`\mathrm{}\omega _z\delta _{n_xn_x^{}}\delta _{n_yn_y^{}}\delta _{n_zn_z^{}+2}\text{ }\sqrt{n_z^{}(n_z+1)}`$
$`\mathrm{}\omega _z\delta _{n_xn_x^{}}\delta _{n_yn_y^{}}\delta _{n_zn_z^{}2}\text{ }\sqrt{n_z(n_z^{}+1)}`$
$`+cyclicpermutations]`$ (37)
## 6 Symmetry properties of the nuclear surface
The three surfaces can be written as:
$$\pi =\frac{x^2}{A^2}+\frac{y^2}{B^2}+\frac{z^2}{C^2}1=0$$
(38)
This implies the following properties:
$$\pi (x,y,z)=\pi (x,y,z)=\pi (x,y,z)=\pi (x,y,z)$$
(39)
Thus, for the two mean fields (i.e. central and spin-orbit fields), and for the Coulomb potential, we obtain:
$$V(x,y,z)=V(x,y,z)$$
(40)
$$S(x,y,z)=S(x,y,z)$$
(41)
$$\mathrm{\Phi }^C(x,y,z)=\mathrm{\Phi }^C(x,y,z)$$
(42)
### 6.1 Parity
Because of the relations (40), (41), and (42) the parity is a good quantum number, and the initial matrix decays into two sub-matrices according to the number
$$p_a=(1)^{n_x+n_y+n_z}=\pm 1$$
(43)
Obviously, if $`(n_x+n_y+n_z)`$ is even or odd the parity is respectively positive or negative.
### 6.2 Signature
Furthermore, the Kramers degeneracy is expressed here, by the fact that the eigenvalues are doubly degenerated relatively to the signature quantum number $`q_K`$ , which is defined by:
$$q_K=(1)^{n_x+n_y}\mathrm{\Sigma }=\pm 1/2$$
(44)
Consequently, the secular matrix splits into two sub-matrices, and only one must be considered. The two matrices contain the same set of eigenvalues, but the eigenfunctions are time-reversed each other.
### 6.3 Consequences of these symmetries
The computer code carries out calculations only for one kind of particles. Therefore, in order to take into account both neutrons and protons, the code must be run twice.
Since the Hamiltonian connects only states with the same parity, the computer code is built in such way that it separates the two types of parity $`p_a=\pm 1`$ and performs the calculations separately for them. Consequently, the representative matrix of the Hamiltonian splits into two blocks with a definite parity for each block. The diagonalization is then carried out in each block.
Furthermore, the Kramers degeneracy involves the same eigenvalues for states which are time-reversed each other. For each block of a definite parity, the eigenenergies can be separated into two sets defined by the signature $`q_K=\pm 1/2`$. The code will make calculations only for $`q_K=+1/2`$. The second block $`q_K=1/2`$ will be implicit, and will contain same energies but with time-reversed eigenfunctions. These eigenvectors may be obtained by application of the time reversal operator, i.e. by the operator $`T=i\sigma _yK_0`$ , where $`\sigma _y`$ is a Pauli matrix and $`K_0`$ is the operator of complex conjugation.
Thus, one obtains 8 blocks, of which 4 are actually calculated (i.e. here the four first).
$`1)`$ $`\left[n\right]`$ $`\left[p_a=+1\right]`$ $`\left[q_K=+1/2\right]`$
$`2)`$ $`\left[n\right]`$ $`\left[p_a=1\right]`$ $`\left[q_K=+1/2\right]`$
$`3)`$ $`\left[p\right]`$ $`\left[p_a=+1\right]`$ $`\left[q_K=+1/2\right]`$
$`4)`$ $`\left[p\right]`$ $`\left[p_a=1\right]`$ $`\left[q_K=+1/2\right]`$
$`5)`$ $`\left[n\right]`$ $`\left[p_a=+1\right]`$ $`\left[q_K=1/2\right]`$
$`6)`$ $`\left[n\right]`$ $`\left[p_a=1\right]`$ $`\left[q_K=1/2\right]`$
$`7)`$ $`\left[p\right]`$ $`\left[p_a=+1\right]`$ $`\left[q_K=1/2\right]`$
$`8)`$ $`\left[p\right]`$ $`\left[p_a=1\right]`$ $`\left[q_K=1/2\right]`$
So, it is important to point out that, the number of the basis states practically taken into account by the code is the half of the actual number.
## 7 Numerical choices and prescriptions
### 7.1 The quadratures
The matrix elements of $`V(\stackrel{}{r})`$, $`e\mathrm{\Phi }^C(\stackrel{}{r})`$ and $`V^{so}(\stackrel{}{r})`$ are calculated with the Gauss-Hermite method, with $`30\times 30\times 30`$ of mesh points. The Coulomb potential $`\mathrm{\Phi }^C(\stackrel{}{r})`$ is also evaluated numerically by the Gauss-Legendre method, but with $`48\times 48`$ of mesh points .
These choices seem to be sufficient relatively to the size of the basis ( $`N_{\mathrm{max}}26`$ ), and the interval of deformation ( $`0\beta 0.6`$). A direct checking has been done by increasing the number of quadrature points and by comparing the stability of the results (even with 20 points the results remain very correct).
### 7.2 Prescription of the basis truncation
In practice, the Hamiltonian matrix is finite. Therefore, for reasons of accuracy, we have to select a sufficient number of the basis states. Generally, we adopt one of the two following criteria:
The first (spherical criterion) consists in choosing all basis states which satisfy the following inequality.
$$n_x+n_y+n_zN_{\mathrm{max}}$$
(45)
With this criterion the total number of the basis states is given by (Nmax+1) x (Nmax+2) x (Nmax+3)/6.
In the second criterion (deformed criterion), one selects the states according to the deformation of the basis, i.e. according to the three frequencies of basis.
$$(n_x+\frac{1}{2})\mathrm{}\omega _x+(n_y+\frac{1}{2})\mathrm{}\omega _y+(n_z+\frac{1}{2})\mathrm{}\omega _zE_{cut}=(N_{\mathrm{max}}+\frac{3}{2})\mathrm{}\omega _0$$
(46)
(In fact, these three frequencies are already connected by the condition (27)).
Thus, the choice of $`N_{\mathrm{max}}`$ determines the size of the basis. The files ”conver12.res” and ”conver13.res” give some details about this.
### 7.3 Optimization of the basis frequencies
Since the Hamiltonian operator does not depend on the oscillator frequencies, its eigenfunctions, and its eigenenergies, must not depend on these parameters. In practice, the representative matrix of the Hamiltonian is built by means of a finite number of oscillator eigenfunctions. This implies a *spurious* dependence according to these parameters.
In another point of view, we might consider this method as a variational method in which the variational parameters are the frequencies of the basis. Thus, the best set (in terms of energy) for these frequencies should be precisely the one, *which minimizes the eigenenergies*, or simply their sum.
For practical reasons, this method is not easy, since the variation is three-dimensional. However, it can be often more efficient to use some prescriptions in order to find (in an economical way) suitable values for these parameters.
In the present work, we have adopted the approach of the references and . In that method, we define first, the quantities $`p`$ and $`q`$ by:
$$q^2=\frac{z^2}{x^2}=\frac{𝑑\tau .\rho (\stackrel{}{r})z^2}{𝑑\tau .\rho (\stackrel{}{r})x^2}p^2=\frac{z^2}{y^2}=\frac{𝑑\tau .\rho (\stackrel{}{r})z^2}{𝑑\tau .\rho (\stackrel{}{r})y^2}$$
(47)
where $`\rho (\stackrel{}{r})`$ is the nuclear density. Note that the present definition of $`p`$ differs from that of the ref..
For a harmonic oscillator, the equations (47) are reduced to very simple relations.
$$q_{HO}=\frac{\omega _x}{\omega _z}p_{HO}=\frac{\omega _y}{\omega _z}$$
(48)
Next, we have to add, to these two formulas, the relation (27). Now, it is possible to replace the parameter set ( $`\omega _x,\omega _y,\omega _z`$ ) by the equivalent ($`q_{HO},p_{HO},\omega _0`$).
In the same way, for the potential of Woods-Saxon, the nuclear density can be approximated by the one of the liquid drop (i.e. a constant density). We obtain thus:
$$q_{WS}=\frac{c}{a}p_{WS}=\frac{c}{b}$$
(49)
At last, we “adapt” the oscillator basis to the nuclear shape by requiring:
$$q_{HO}=q_{WS}p_{HO}=p_{WS}$$
(50)
For the $`\omega _0`$ value, we can adopt simply the one of the Nilsson model.
$$\mathrm{}\omega _041.A^{.\frac{1}{3}}$$
(51)
Many tests have shown that relations (50) and (51) give *automatically* very close values to those that produce the ”true” minimization. Furthermore, a general rule is that a large basis size involves always a weak dependence of the eigenvalues according to these parameters. Going to the limit, we can say that if the basis was infinite, the results would be independent to the basis parameters. Conversely, for a too small basis, the dependance is strong, and the results become too inaccurate.
For a square well, or (approximately) a Woods-Saxon potential, simple analytical considerations lead to a more ”refined” value for the parameter $`\mathrm{}\omega _0`$:
$$\mathrm{}\omega _0\frac{5}{3}\left(\frac{2}{\pi ^2}\right)^{1/3}\left|V_0\right|.A^{1/3}0.979\left|V_0\right|.A^{1/3}$$
(52)
where $`V_0`$ is the depth of the potential. The equation (52) is obtained by requiring the condition
$$r^2_{harm.\text{ }Oscillator}=r^2_{square\text{ }well}$$
The averages are made with the semi-classical Thomas-Fermi density:
$$\rho _{TF}(r)=\frac{(2m)^{3/2}}{3\pi ^2\mathrm{}^3}(\lambda V(r))^{3/2}$$
The Fermi level $`\lambda `$ is determined by the condition of conservation of the particle number $`A`$.
The relation (52) could explain the empirical scale factor used sometimes in the ” standard equation” (51).
### 7.4 The numerical values of the $`\beta `$ parameters of the basis
The quantities $`\beta _0=\sqrt{\frac{m\omega _0}{\mathrm{}}}`$ , $`\beta _x=\sqrt{\frac{m\omega _x}{\mathrm{}}}`$ , $`\beta _y=\sqrt{\frac{m\omega _y}{\mathrm{}}}`$ , and $`\beta _z=\sqrt{\frac{m\omega _z}{\mathrm{}}}`$. are numerically calculated like
$$\sqrt{\frac{m\omega }{\mathrm{}}}=\sqrt{\frac{mc^2}{\mathrm{}^2c^2}\mathrm{}\omega }$$
The values are:
$`m_pc^2=938.2592MeV`$
$`m_nc^2=939.553MeV`$
$`\mathrm{}c=197.32879MeVfm`$
This involves:
$$\frac{m_pc^2}{(\mathrm{}c)^2}=t_p=0.0240958315MeV^1fm^2$$
(53)
$$\frac{m_nc^2}{(\mathrm{}c)^2}=t_n=0.0241290571MeV^1fm^2$$
(54)
so that :
$$\beta _0=\sqrt{t.\mathrm{}\omega _0},\beta _x=\sqrt{t.\mathrm{}\omega _x},\beta _y=\sqrt{t.\mathrm{}\omega _y},\beta _z=\sqrt{t.\mathrm{}\omega _z}$$
(55)
with $`t=t_p`$ or $`t=t_n`$
The numbers $`t_p`$ and $`t_n`$ appear in the subroutine “Basisparam”.
## 8 Diagonalization
The diagonalization of the representative matrix of the Hamiltonian is carried out by a set of subroutines extracted from the EISPACK library of FORTRAN programs (http//www.netlib.org/eispack/). Thus, four subroutines of this library were gathered:
The subroutine tred1 transforms any full symmetrical matrix into a tridiagonal symmetrical matrix by using the Givens-Householder’s method.
For a tridiagonal symmetrical matrix the subroutine tql1 uses the ql method to calculate only the eigenvalues of a tridiagonal matrix.
For a tridiagonal symmetrical matrix, the subroutine tsturm calculates the eigenvalues contained in a given interval. This subroutine calculates also the eigenvectors associated to the found eigenvalues. The adopted method is that of the bisection and the inverse iteration.
Lastly, the subroutine trbak1 recalculates the eigenvectors found by tsturm relatively to the initial basis (that of tred1). The sought eigenvectors are thus obtained.
These subroutines are called by the subroutines diagoplus (for the positive parity), and diagominus (for the negative parity) in which the options of the diagonalization are specified. These options are indicated in the comments of the program, and below, in the subsection 10.1.1.
## 9 The subroutines and the functions of the Program.
The program is composed by a main program, 29 subroutines and 6 functions. The role reserved for each program is briefly described in the paragraph below (and described again in details in the comments of the program).
In fact, all calculations are governed by the subroutine *setsub* which is in some sense a super subroutine.
### 9.1 The set of subroutines (in the order of the calls)
1. The subroutine *read1*: reads the basic input parameters in the file input.dat
2. The subroutine *write1*: performs some tests and writes on a files eigvals.res and conver.res
3. The subroutine *setsub*: drives the successive calculations
4. The subroutine *write2*: writes on the file eigvals.res
5. The subroutine *write3*: writes on the file conver.res
6. The subroutine *woodsparam*: calculates the Myers parameters
7. The subroutine *surfparam*: calculates the surfaces parameters
8. The subroutine *basisparam*: calculates the oscillator basis parameters
9. The subroutine *pottablo* : stores the potential at the nodes of quadrature
10. The subroutine *coefftablo*: stores the products of the coefficient of quadrature
11. The subroutine *hermitablo*: stores the Hermite polynomials at the nodes
12. The subroutine *coultablo*: stores the coulomb potential at the nodes
13. The subroutine *statesplus*: selects the numbers and the oscillator basis states corresponding to the positive parity
14. The subroutine *statesminus*: selects the numbers and the oscillator basis states corresponding to the negative parity
15. The subroutine *idm*: calculates the total numbers of the used basis states
16. The subroutine *matpotplus*: calculates the representative matrix of the central mean potential for the positive parity
17. The subroutine *matpotminus*: calculates the representative matrix of the central mean potential for the negative parity
18. The subroutine *matcinplus*: calculates the representative matrix of the kinetic energy for the positive parity
19. The subroutine *matcinminus*: calculates the representative matrix of the kinetic energy for the negative parity
20. The subroutine *matpotsoplus*: calculates the representative matrix of the mean spin-orbit energy for the positive parity
21. The subroutine *matpotsominus*: calculates the representative matrix of the mean spin-orbit energy for the negative parity
22. The subroutine *diagoplus*: diagonalizes the representative matrix of the hamiltonian for the positive parity
23. The subroutine *diagominus*: diagonalizes the representative matrix of the hamiltonian for the negative parity
24. The subroutine *tred1*: Eispack subroutine (see section 8)
25. The subroutine *tql1*: Eispack subroutine (see section 8)
26. The subroutine *tsturm*: Eispack subroutine (see section 8)
27. The subrouine *trbak1*: Eispack subroutine (see section 8)
28. The subroutine *eigenvalues*: gathers the eigenvalues for both parities
29. the subroutine *vektors*: writes the eigenfunctions in a file.
Fore several subroutines, the names ending in “plus” or in“minus” means that the subroutine performs calculations specifically for a defined parity. The term “plus” is employed for the positive parity, and the term “minus” for the negative parity.
### 9.2 The set of functions
1. The function *Hermite*: calculates the Hermite polynomials
2. The function *delta*: delta symbol of Kroneker
3. The function *potenv*: calculates the central mean potential value at any point.
4. The function *potenso*: calculates the spin-orbit mean potential value at any point.
5. The function *ephi*: calculates the Coulomb energy of the proton at any point.
6. The function *epslon*: estimates the round-off error for the Eispack subroutines
## 10 Input-output data of the FORTRAN program
If no modifications are made the use of the program as presented in long theoretical description is very simple.
### 10.1 The input data
All input data are read from two files in a namelist type declarations. The second file is needed only if one does use a personal parameters for the potential, instead those of Myers. In this latter case, one has to precise its own parameters, in a second separate file.
#### 10.1.1 The first input data file: *input.dat*
The file input.dat gathers all basic input data. Their significance is given below.
* *nmax1* and *nmax2* are the bounds of the loop for $`N_{\mathrm{max}}`$(eq(45) or (46)). This latter is the number of the major shells used in the calculations. If nmax1 =nmax2 (=nmax) the calculations are performed once. The variation of nmax is envisaged only if one desires to study the convergence of the results as a function of the number of the basis states.
* *pi* is the pi number (3.1415927410125d.0)
* If *kkind=1*, calculations are made for the neutrons case.
If *kkind=2*,calculations are made for the protons case.
Any other value of the kkind parameter involves an error declaration of the program.
* *Iz* = number of protons.
* *In* = number of neutrons.
* *Betta*, and *gama* are the usual deformation parameters of Bohr (eq.(18a)-(18c)).
* If *ibase=0* the states of the basis are selected according to the spherical criterion (45).
If *ibase=1*. The states of the base are selected according to the deformed criterion (46) There is not other value for this parameter.
* If *i1i2=1*, the program gives all eigenvalues, without eigenvectors.
If *i1i2=2*,the program gives the eigenvalues included in a given interval \[elow, ehigh\] with the corresponding (orthonormalized) eigenvectors. Any other value of this parameter involves an error declaration of the program.
* *Elow* = lower bound of the selected interval.
* *Ehigh*= higher bound of the selected interval.
( Naturally, if this interval is sufficiently large it will contain all eigenvalues. Consequently all eigenvectors will be also given.)
* If *icalc=0*, the parameters of the Woods-Saxon potential are read from the namelist of the second input file parameters dat.
if *icalc=1*, the Myers parameters are calculated by the subroutine woodsparam.
* If *iscal=1* the basis parameter $`\mathrm{}\omega _0`$ is computed from (52) i.e. from $`\mathrm{}\omega _0=0.979\left|V_0\right|.A^{1/3}`$
If *iscal=2* the basis parameter $`\mathrm{}\omega _0`$ is computed from the relation $`\mathrm{}\omega _0=faktor.A^{1/3}.`$(see eq.(51).
* *faktor*= input parameter of the previous relation
#### 10.1.2 The second input data file: *parameters.dat*
There is an option ( governed by the keyword icalc ) in the first input file which permits to the user to employ its own parameters instead of those of Myers.
The data of the file parameters.dat are :
* *v0neut*= deep of the central part of the potential for the neutrons
* *avneut*= diffuseness of the central part of the potential for the neutrons
* *rvneut*= radius of the central part of the potential for the neutrons
* *capasoneu*= spin-orbit coupling strengh for the neutrons
* *assoneu*= diffuseness of the spin-orbit part of the potential for the neutrons
* *rssoneu*= radius of the spin-orbit part of the potential for the neutrons
* *v0pro*= deep of the central part of the potential for the protons
* *avpro*= diffuseness of the central part of the potential for the protons
* *rvpro*= radius of the central part of the potential for the protons
* *capasopro*= spin-orbit coupling strengh for the protons
* *assopro*= diffuseness of the spin-orbit part of the potential for the protons
* *rssopro*= radius of the spin-orbit part of the potential for the protons
* *rchpro*= radius of the coulomb potential
### 10.2 The output data
The global results can be extracted from the five arrays *evalplus, evalminus, evecplus, evecminus,* and *energies*, in the main program.
The arrays *evalplus* and *evalminus* contain respectively, the eigenvalues for the positive parity and the negative parity . The eigenvalues are classified in an increasing order.
In the same way, the arrays *evecplus* and *evecminus* contain the components of the eigenvectors, in columns, in the same order as that of the eigenvalues.
For the positive parity ( respectively the negative parity) , the parameter *nevalplus* in the subroutine diagoplus (respectively *nevalminus* in the subroutine diagominus) gives the number of eigenvalues.
Sometimes, it is more convenient to gather all eigenvalues in a common array (but the eigenvectors remain in their respective blocks). This is carried out in a common array named *energies*. In this array, the eigenvalues are classified in an increasing order.
In order to find the corresponding eigenvector to a given eigenvalue, a vector containing a supplemental information was created and named *num(k)*. The sign and the absolute value of num(k) indicate respectively the block (i.e. evecplus or evecminus) and the place of the column in this block.
Furthermore, the output data can be consulted in a straightforward way, in three files:
a) The eigenvalues are written in the file ” eigvals.res ”. In this file, it is indicated in particular, if the eigenvalues belong to the set corresponding to the positive parity or those corresponding of the negative parity.
b) The eigenfunctions are recorded in the file ” vekt.res ”. For every eigenvalue there is a set of components relative to the different states (nx, ny, nz, sigma) of the basis.
c) A brief study on the convergence is made in the file ” conver.res ”.
## 11 Checking of the computer code and comments on the test run
In order to check the code, one has proceeded to three types of tests. In the first, we use well-known analytical results. In the second, we compare our calculations with those using the same model. At last, in the third, we use some well-known properties of symmetry.
### 11.1 Analytical tests
In fact, the method of resolution of the Schrodinger equation proposed here is a *purely numerical method*. Consequently, one can use, not only the Woods-Saxon potential, but also *any other type of potential*. It is then possible to replace the Woods-Saxon potential by that of the harmonic oscillator in order to test the code by well-known analytical results.
#### 11.1.1 The deformed case without spin-orbit term
Indeed, for a pure deformed harmonic oscillator, (without spin-orbit interaction), in cartesian coordinates, the theoretical expression of the energy is given simply by:
$`E(n_x,n_y,n_z)=(n_x+\frac{1}{2})\mathrm{}\omega _x+(n_y+\frac{1}{2})\mathrm{}\omega _y+(n_z+\frac{1}{2})\mathrm{}\omega _z`$
$`n_x,n_y,n_z=0,1,2,\mathrm{}\mathrm{}.\mathrm{}`$
For reasons of simplicity, we have chosen the same frequencies as those of the basis. The numerical values are extracted from the file ”eigvals1.res”.
In the computer program, one must replace the Woods-Saxon potential by that of the harmonic oscillator, i.e. by:
$`V(\stackrel{}{r})=\frac{1}{2}m(\omega _x^2x^2+\omega _y^2y^2+\omega _z^2z^2)`$
Then, one has to cancel the spin-orbit interaction (by making the function potenso = 0 or by cancelling the spin-orbit coupling constant) in the code.
Calculating some levels analytically, and comparing them with those of the code, one can note an excellent agreement (to seven significant digits) for the deformed case (see table 1)
It is important to note that the matrix elements are integrated numerically in the computer code, therefore, from this test, we can conclude that the program performs this task correctly. Because this term is diagonal ( in fact the code ”does not know this” but after calculations, it finds that the nondiagonal elements are equal to zero ) in our basis, this test does not permit us to verify the diagonalization. These latter part of the program will be verified in the subsections below.
Now, if we add the spin-orbit interaction, we could test the program entirely. Unfortunately, in the deformed case, there is no theoretical expression for that.
#### 11.1.2 The spherical case without, and, with spin-orbit term
Of course, for the spherical case, it is possible to make $`\mathrm{}\omega _x=\mathrm{}\omega _y=\mathrm{}\omega _z=\mathrm{}\omega _0`$ in the previous theoretical expression. Nevertheless, the spherical coordinates are more convenient because as we shall see, the spin-orbit term has to be ”treated” in that system. In this latter, the theoretical expression of the energy of a pure oscillator is well-known:
$`E(n,l,m)=\left[2(n1)+\mathrm{}+\frac{3}{2}\right]\mathrm{}\omega _0=\left[N+\frac{3}{2}\right]\mathrm{}\omega _0`$
$`N=2(n1)+\mathrm{}=`$ This number specifies a major shell
$`n=1,2,\mathrm{}\mathrm{}\mathrm{}.\mathrm{}`$
$`\mathrm{}=0,1,2,\mathrm{}\mathrm{}\mathrm{}\mathrm{}`$
$`m=\mathrm{},\mathrm{}+1,\mathrm{},\mathrm{}`$, (for a fixed $`\mathrm{}`$)
Due to the fact that the spherical symmetry is a particular case of the deformed case, it is obvious, that the results of the code (eivals2.res ) should be in complete agreement ( with the awaited degeneracy) with the analytical results. We can see in.table 2a,.that it is really the case.
Furthermore, there, contrary to the deformed case, it is possible to obtain the analytical expression for the spin-orbit term.
Indeed, one can use the relation (22) in a suitable way in order to obtain a simple theoretical expression for the spin-orbit term.
Taking $`S(\stackrel{}{r})=cr^2/2`$, where c is a positive constant, one gets then:
$`\frac{1}{r}\frac{S(r)}{r}\stackrel{}{\mathrm{}}.\stackrel{}{\sigma }=c\stackrel{}{\mathrm{}}.\stackrel{}{\sigma }=2c\stackrel{}{\mathrm{}}.\stackrel{}{s}=2c\frac{1}{2}(\stackrel{}{j}^2\stackrel{}{\mathrm{}}^2\stackrel{}{s}^2)`$
The most important point is that, in this way the splitting of the major shells does not depend on $`r`$, and, is then *rigorously* given by:
$`\mathrm{\Delta }E(\mathrm{}\frac{1}{2})=c.(\mathrm{}+1)`$ if $`j=\left|\mathrm{}\frac{1}{2}\right|`$ and $`\mathrm{}0`$
$`\mathrm{\Delta }E(\mathrm{}+\frac{1}{2})=c.\mathrm{}`$ if $`j=\mathrm{}+\frac{1}{2}`$ and $`\mathrm{}0`$
Thus, the new energies can be written as:
$`E(n,l,j=\mathrm{}\frac{1}{2})=\left[2(n1)+\mathrm{}+\frac{3}{2}\right]\mathrm{}\omega _0+\mathrm{\Delta }E(\mathrm{}\frac{1}{2})`$
Therefore, the code can be verified in its integrality.
In order to simplify the numerical values of the splitting, we take $`c=1MeV`$. Thus, except for the value $`\mathrm{}=0,`$ we can see that the levels are simply shifted by integer values according to the value of $`\mathrm{}`$. In order to illustrate that, we will give two examples:
* Example 1 :
if $`\mathrm{}=1`$, the $`p`$ shell with a energy noted $`E(p)`$ splits into two subshells according to the two values of $`j`$:
for $`j=\mathrm{}+\frac{1}{2}=1+\frac{1}{2}=3/2`$,
$`E(p3/2)=E(p)\mathrm{}=E(p)1Mev`$
for $`j=\mathrm{}\frac{1}{2}=1\frac{1}{2}=1/2`$,
$`E(p1/2)=E(p)+(\mathrm{}+1)=E(p)+2Mev`$
* Example 2:
similarly, if $`\mathrm{}=3`$ ($`f`$ shell ) one obtains
for $`j=\mathrm{}+\frac{1}{2}=3+\frac{1}{2}=7/2`$,
$`E(f7/2)=E(f)\mathrm{}=E(f)3Mev`$,
for $`j=\mathrm{}\frac{1}{2}=3\frac{1}{2}=5/2`$,
$`E(f5/2)=E(f)+(\mathrm{}+1)=E(f)+4Mev`$
In the table 2b, we compare all results of the code (eigvals3.res) with those of the analytical expression. Practically, the code ( which works in double precision) gives the exact values to six or seven significant digits for all levels of the spectrum.
This high accuracy is due to the fact that the oscillator potential is a polynome of order two, therefore the Gauss method gives in this case the exact values for all matrix elements.
Of course, these tests are not realistic cases, but they prove that the code runs properly with a high degree of precision.
The case with spin-orbit term is very important because it involves the integral analytical checking of the code. Due to the fact that this operator is not diagonal in the oscillator basis, it proves not only that the code performs correctly all calculations of the matrix elements, but also proves that the step of the diagonalization is done properly.
One also made some additional easy checks (not shown here). For example, by taking a constant potential in the spin-orbit term , one cancels the spin-orbit potential . That was well verified by the code, etc…
### 11.2 Comparisons with similar works
For the deformed Woods-Saxon potential , it seemed to us more convenient to compare our code with those of the reference . The reasons are the following:
a) We use exactly the same model as this reference.
b) All potential parameters of the calculations are precised in that reference, and we need to use the same.
c) Not only a part, but the entire spectrum of eigenvalues is given (as a function of the deformation).
The only disadvantage is that the results are displayed under a graphical form. However, in extracting the numerical values, we have tried to minimize the errors by using a graphical software.
The eigenvalues are read with the own scale of the software. Then, a suitable linéar transformation returns these values in MeV. Nevertheless, in order to find the ”best values”, this transformation has been carried out by the least-squares’ method of the software.
It turns out that it is possible to obtain values with an error about $`\pm 0.03MeV.`$
We have thus considered the deformation $`(\beta =0.3,\gamma =0.0)`$ for the lead Pb208.
For the basis parameters ($`Nmax`$ and $`\mathrm{}\omega _0`$), we tried to use in calculations, the same, in order to obtain, as much as possible, close results. For Nmax, the reference indicates that the matrices corresponding to the two parities have a dimension of about 160 states. Consequently the fixed value for Nmax was certainly Nmax=10. However, the value of $`\mathrm{}\omega _0`$ really used by the code is not given. This reference indicates only that, for the spherical case, the theoretical relation $`\mathrm{}\omega _0=55.A^{1/3}`$ is better than the standard theoretical relation $`\mathrm{}\omega _0=41.A^{1/3}`$. Nevertheless, the reference , claims that a practical value of the order of 45-48MeV (instead $`55.MeV`$) gives a somewhat better results that these theoretical relations. Since the codes of these two references have been compared, it is probable that a common practical value was fixed. We endeavored ” to guess ” this value. After many tests, It turned out that the value 47 MeV gave a good agreement
Our calculations were carried out successively with $`Nmax=10`$ (as the cited reference), and $`Nmax=26`$. Indeed, this latter value insures that the levels are calculated with about three or four significant digits near the fermi level, and obviously, all the more for lower levels (see the file ”conver13.res”). They are thus practically independent of the choice of the basis parameters
In the tables (3a-3b,4a-4b) which have been deduced from the files ”eigvals4.res”, ”eigvals5.res”, ”eigvals6.res”, ”eigvals7.res”, we show respectively all bound levels of the Pb208 for four cases:
a) *neutrons-prolate shape (*$`\gamma `$*=0*$`{}_{}{}^{},\beta =0.3`$*),Nmax=10*
b) *protons-prolate shape (*$`\gamma `$*=0*$`{}_{}{}^{},\beta =0.3`$*),Nmax=10*
c)*neutrons-prolate shape (*$`\gamma `$*=0*$`{}_{}{}^{},\beta =0.3`$*)Nmax=26*
d)*protons-prolate shape (*$`\gamma `$*=0*$`{}_{}{}^{},\beta =0.3`$*)Nmax=26*
The levels were separated in two distinct blocks according to their parity.
Of course, for a finite potential, the discrete positive energy levels do not represent, a valid solution of the continuum (see ref.), therefore, we shall drop them .
In all cases, we can note that the energy levels are practically the same ones for the low part of the spectrum, but relative small differences appear in the *upper part of the spectrum*.
These differences are more prounonced for $`Nmax=26`$ that for $`Nmax=10`$. The analyse of these results leads to the following conclusions:
* The lowest levels of the spectrum converge systematically more quickly than the others. As one goes up in the spectrum the convergence is in general slowest, but there can be some rare exceptions.
* A rapid convergence involves a weak dependence relatively to the basis parameters. The highest levels of the spectrum are thus more sensitive to the basis parameters. One can affirm that if the basis parameters of our code are close to those of the reference , they are not rigorously the same ones.
* In fact, one noted that this remark is general. Indeed, a modification of any parameter (for example those of the potential) in the calculations produces a modification relatively more significant for the highest levels than for the lowest levels. For example if the radius of the mean potential (spherical case for the neutrons) varies from 7.36fm to 7.40fm (all other parameters being constant), the first level, and the Fermi level undergo variations of 0.05MeV, and 0.28MeV respectively. The ”general rule” is thus that *the lowest levels are most ”stable”*.
* Owing to the fact that we employ very similar parameters, our results with $`Nmax=10`$ are ”artificially” very close to those of the reference (the mean deviations are about 0.05 MeV for all cases). Thus, our purpose which was to recover the same results is now reached. But, the word ”artificially” means that for this small basis, both results are not enough accurate, although they are the same
Indeed, it is clear that they will be actually less precise than those obtained with $`Nmax=26`$. Significant differences appear in the top of the spectrum. In the file ”conver13.res”, one can note that the Fermi level is stabilized to about $`0.010.03`$ $`MeV`$ only starting from $`Nmax=1516`$. Therefore, calculations with $`Nmax=1014`$ produce mediocre results.
The rapid convergence of the lowest states is due mainly to the fact that the corresponding wave functions are very similar to those of the oscillator. This is not the case for the highest states where the wave functions are strongly oscillating, and where the edge effect of the potential is ”felt”.
This can be easily noticed in the components of the eigenvectors, in the file ”vekt14.res” . For example, concerning the first eigenvalue, only the components corresponding to the lowest quantum numbers are important (see the components numbered 1, 2, 12, and, 59).
### 11.3 Tests using some properties of the parametrisation $`(\beta ,\gamma )`$
Two simple tests can be carried out to check the consistence of the program:
In the first, one compares the spectra obtained with the deformations$`(\beta ,\gamma )`$ and $`(\beta ,\gamma )`$ This operation is in fact nothing other that a simple permutation of the axes 1 and 2 of the ellipsoid. Of course the two shapes are the same, consequently, the respective spectra must be identical.
In the files ”eigvals8.res” and ”eigvals9.res” one can easily check that is really the case with an astonishing precision. In particular, one can note in these files the permutation of values of the parameters $`\mathrm{}\omega _x`$(hbaromegx), and $`\mathrm{}\omega _y`$(hbaromegy).
In the second, one compares the spectra obtained with the deformations $`(\beta ,\gamma =60{}_{}{}^{})`$ and $`(\beta ,\gamma =0{}_{}{}^{})`$. There also, this operation is simply a cyclic permutation of the three axes of the ellipsoid. Therefore, the spectrum must also remain unchanged.
As for the previous case, this can be easily verified in the files ”eigvals10.res and eigvals11.res”.
### 11.4 Tests of convergence
In the files ”conver12.res” and ”conver13.res”, we have shown the convergence of the sum (of the single particle energy) of the first 126 neutrons levels of Pb208 for two deformations. The potential’s parameters are those of the reference .
This sum has converged to less than 1 Mev only starting from the values $`Nmax=14`$ and $`Nmax=16`$, respectively for the spherical and the deformed cases.
This implies for the Fermi level, a convergence to 0.02 MEV and 0.01 MEV respectively for these two cases. However, for $`Nmax16`$, theses deviations depend still of the value $`\mathrm{}\omega _0`$. Obviously, for higher bound states , the precision will be less.
Everything depends on what one wants to make. So, for example , for the Strutinsky’s shells corrrection the previous values seem to be sufficient.
Always concerning the Fermi’s level (conver13.res), one notices in general that it increases in absolute value as $`Nmax`$ increases, but sometimes, it happens that it decreases slightly (in absolute value). For example, in the spherical case, it passes from 8.510 to 8.502 when Nmax passes from 13 to 14. We can easily see that the dimensions of subspaces corresponding to the positive and negative parity do not vary simultaneously when Nmax varies by one unit. For example, when one passes from Nmax=13 to Nmax=14 only one subspace, namely the one with a positive parity, undergoes changes from 252 basis states to 372 basis states. The other remains the same with 308 basis states.
In our example, the Fermi level belongs to the subspace of negative parity, therefore, apparently, it should not have to change. In fact, the formulae of the spin-orbit interaction connects the matrix elements of the two subspaces (see eq.34-36). This implies always a slight modifications in the subspace which has not varied, and this must not be assimilated to a noise.
In fact, in this method, the ”true noise” has two main sources :
a) under-estimations of the number of points in the numerical integrations of the matrix elements.
b) a too small basis or really inadequate values of the basis parameters.
With 30 points of quadrature, a double precision, and a large basis ($`Nmax`$ up to 26) these two problems are here minimized.
## 12 Conclusion
We have elaborated and checked a calculation program solving the equation of Schrodinger for a deformed potential of Woods-Saxon type.
The program appears very rapid, and consequently, it becomes possible to use significant basis sizes.
Calculations with small bases, like those which were carried out in the past with $`Nmax=1012`$ lead to a very poor precision. Our conclusions are corroborated by other works. For example, the ref. has shown for Hartree-Fock calculations that the error in the energy of Pb208 is smaller than 1MeV only for $`Nmax16`$. Other examples are given in the ref. which confirm this fact.
Similar codes were made in the past, but with the assumption of axial symmetry. To our knowledge, triaxial Woods-Saxon calculations were never really undertaken with significant sizes of the oscillator basis.
## Appendix A The Myers parameters
The diffuseness parameters $`a_V`$, $`a_{so}`$, and the spin-orbit coupling $`\kappa `$ are the same as those of the Ref..
$$a_V=0.66fm$$
(56)
$$a_{so}=0.55fm$$
(57)
$$\kappa =12MeVfm^2$$
(58)
The parameters of central potential, and of the spin-orbit potential, were extracted from the Myers droplet model . This theory uses Thomas Fermi’s approximation to approach average properties of finite nuclei like the density radii, skin-thicknesses, …, in terms of neutron and proton numbers.
In this model, two auxiliary quantities are first defined:
$$\widehat{\delta }=\frac{\frac{NZ}{A}+0.0112\frac{Z^2}{A^{5/3}}}{1+\frac{3.15}{A^{1/3}}}$$
(59)
$$\widehat{ϵ}=\frac{0.147}{A^{1/3}}+0.330\widehat{\delta }^2+\frac{0.00248Z^2}{A^{4/3}}$$
(60)
The physical significance of these quantities is explained in the Ref.
With help of these quantities, the depth of the mean potentials are written as:
$$V_0(protons)=52.548.7\widehat{\delta }$$
(61)
$$V_0(neutrons)=52.5+48.7\widehat{\delta }$$
(62)
The radii of the central potentials (which are different for protons and neutrons) are expressed by means of the nuclear density radii $`R_0(protons)`$, or, $`R_0(neutrons)`$, and the density diffuseness $`a_V`$:
$$R_V(protons)=R_0(protons)\left\{1\frac{\pi ^2}{3}\left(\frac{a_V}{R_0(protons)}\right)^2\right\}$$
(63)
$$R_V(neutrons)=R_0(neutrons)\left\{1\frac{\pi ^2}{3}\left(\frac{a_V}{R_0(neutrons)}\right)^2\right\}$$
(64)
with
$$R_0(protons)=R_0+0.82\frac{0.56}{R_0}+0.22\widehat{\delta }$$
(65)
$$R_0(neutrons)=R_0+0.82\frac{0.56}{R_0}0.22\widehat{\delta }$$
(66)
$$R_0=r_0A^{1/3}(1\widehat{ϵ})$$
(67)
$$r_0=1.16fm$$
(68)
The radius of the spin-orbits mean field is given in the same way:
$$R_{so}=R_0\left(1\frac{\pi ^2}{3}\left(\frac{a_{so}}{R_0}\right)^2\right)$$
(69)
At last, the radius of the charge density is given by:
$$R_{ch}=R_0\frac{1}{3}r_0A^{1/3}\left(\frac{NZ}{N+Z}\widehat{\delta }\right)$$
(70)
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# Strong-coupling theory of high-temperature superconductivity and colossal magnetoresistance
## I INTRODUCTION: The “Fröhlich-Coulomb” model
Although high-temperature superconductivity (HTS) has not yet been targeted as ‘the shame and despair of theoretical physics’, - a label attributed to low-temperature superconductivity during the first half-century after its discovery - controversy of current theoretical constructions has led many researchers to say that there is no theory of HTS and no progress in understanding the phenomenon. A significant fraction of theoretical research in the field has suggested that the interaction in novel superconductors is essentially repulsive and unretarted, and it could provide high $`T_c`$ without phonons. Indeed strong onsite repulsive correlations (Hubbard $`U`$) are essential in shaping the insulating state of undoped (parent) compounds. Different from conventional band-structure insulators with completely filled and empty Bloch bands, the Mott insulator arises from a potentially metallic half-filled band as a result of the Coulomb blockade of electron tunnelling to neighboring sites mott . However, the Hubbard $`U`$ model shares an inherent difficulty in determining the order when the Mott-Hubbard insulator is doped. While some groups have claimed that it describes high-$`T_c`$ superconductivity at finite doping, other authors could not find any superconducting instability. Therefore it has been concluded that models of this kind are highly conflicting and confuse the issue by exaggerating the magnetism rather than clarifying it lau .
The Hubbard-$`U`$ model of high temperature superconductivity or its strong-coupling $`\mathrm{"}tJ\mathrm{"}`$ approximation are also refutable on the experimental ground. The characteristic magnetic interaction, which is allegedly responsible for the pairing in the model, is the spin exchange interaction, $`J=4t^2/U`$, of the order of $`0.1`$ eV (here $`t`$ is the hopping integral). On the other hand, a simple parameter-free estimate of the Fröhlich electron-phonon interaction (routinely neglected within the Hubbard $`U`$ approach) yields the effective attraction as high as $`1`$ eV alebook . This estimate is obtained using the familiar expression for the polaron level shift, $`E_p,`$ the high-frequency, $`ϵ_{\mathrm{}}`$, and the static, $`ϵ_0,`$ dielectric constants of the host insulator, measured experimentally alebra ,
$$E_p=\frac{1}{2\kappa }_{BZ}\frac{d^3q}{\left(2\pi \right)^3}\frac{4\pi e^2}{q^2},$$
(1)
where $`\kappa ^1=ϵ_{\mathrm{}}^1ϵ_0^1`$ and the size of the integration region is the Brillouin zone (BZ). Since $`ϵ_{\mathrm{}}=5`$, $`ϵ_0=30`$ in La<sub>2</sub>CuO<sub>4</sub> and $`ϵ_{\mathrm{}}=3.9`$, $`ϵ_0=16`$ in LaMnO<sub>3</sub> one obtains $`E_p=0.65`$ eV and $`E_p=0.88`$ eV in La<sub>2</sub>CuO<sub>4</sub> and LaMnO<sub>3</sub>, respectively. Hence the attraction, which is about $`2E_p`$, induced by the lattice deformation in cuprates and manganites is one order of magnitude larger than the exchange (magnetic) interaction. There is virtually no screening of e-ph interactions with $`c`$axis polarized optical phonons in cuprates because the upper limit for the out-of-plane plasmon frequency ($`<200`$ cm<sup>-1</sup>)mar is well below the characteristic phonon frequency, $`\omega `$ 400 - 1000 cm <sup>-1</sup> . The screening in manganites is also very poor since the mobility of carriers is very low. As a result of poor screening the magnetic interaction remains small compared with the Fröhlich interaction at any doping. Further compelling evidence for the strong e-ph interactions has come from the isotope effects in cuprates ZHAO and manganites zhao1 , recent high resolution angle resolved photoemission spectroscopies LAN , and a number of earlier optical mic ; ita and neutron-scattering ega studies. Hence any realistic approach to HTS in cuprates, other doped oxides and fullerenes, and to CMR in ferromagnetic oxides should treat the long-range Coulomb and *unscreened* e-ph interactions on an equal footing.
In the past decade we have developed a ”Fröhlich-Coulomb” model (FCM) ALEXAND ; alekor ; alebook to deal with the strong long-range Coulomb and e-ph interactions in cuprates, manganites and other related compounds. The model Hamiltonian explicitly includes a long-range electron-phonon and the Coulomb interactions as well as the kinetic and deformation energies. The implicitly present large Hubbard $`U`$ term prohibits double occupancy and removes the need to distinguish fermionic spins since the exchange interaction is negligible compared with the direct Coulomb and the electron-phonon interactions. The model also provides a simple explanation of CMR in ferromagnetic oxides if the exchange interaction of p-holes with d-electron spins is included in the Hamiltonian alebra2 (see below). Introducing spinless fermionic, $`c_𝐧`$, and phononic, $`d_{𝐦\alpha }`$, operators the Hamiltonian of the model is written as
$`H=`$ $``$ $`{\displaystyle \underset{𝐧𝐧^{}}{}}\left[t\left(𝐧𝐧^{}\right)c_𝐧^{}c_𝐧^{}V_c\left(𝐧𝐧^{}\right)c_𝐧^{}c_𝐧c_𝐧^{}^{}c_𝐧^{}\right]`$ (2)
$``$ $`{\displaystyle \underset{\mathrm{𝐧𝐦}}{}}\omega _\alpha g_\alpha \left(𝐦𝐧\right)\left(𝐞_\alpha 𝐮_{𝐦𝐧}\right)c_𝐧^{}c_𝐧\left(d_{𝐦\alpha }^{}+d_{𝐦\alpha }\right)`$
$`+`$ $`{\displaystyle \underset{𝐦\alpha }{}}\omega _\alpha \left(d_{𝐦\alpha }^{}d_{𝐦\alpha }+1/2\right),`$
where $`𝐞_\alpha `$ is the polarization vector of the $`\alpha `$th vibration coordinate, $`𝐮_{𝐦𝐧}\left(𝐦𝐧\right)/\left|𝐦𝐧\right|`$ is the unit vector in the direction from electron $`𝐧`$ to ion $`𝐦`$, $`g_\alpha \left(𝐦𝐧\right)`$ is the dimensionless e-ph coupling function, and $`V_c\left(𝐧𝐧^{}\right)`$ is the inter-site Coulomb repulsion. $`g_\alpha \left(𝐦𝐧\right)`$ is proportional to the force acting between the electron on site $`𝐧`$ and the ion on $`𝐦`$. For simplicity, we assume that all the phonon modes are non-dispersive with the frequency $`\omega _\alpha `$. We also use $`\mathrm{}=k_B=c=1`$.
The Hamiltonian, Eq.(2), has been solved analytically by using the $`\mathrm{"}1/\lambda \mathrm{"}`$ multi-polaron expansion technique alebook in the strong limit where the e-ph coupling constant is large, $`\lambda =E_p/zt>1`$. Here the polaron level shift is $`E_p=_{𝐧\alpha }\omega _\alpha g_\alpha ^2\left(𝐧\right)\left(𝐞_\alpha 𝐮_𝐧\right)^2,`$ and $`zt`$ is a half-bandwidth in the rigid lattice. The model shows a reach phase diagram depending on the ratio of the inter-site Coulomb repulsion $`V_c`$ and the polaron level shift $`E_p`$ alekor . The ground state of FCM is a *polaronic* Fermi liquid when the Coulomb repulsion is large, a *bipolaronic* high-temperature superconductor at intermediate Coulomb repulsions, and a charge-segregated insulator if the repulsion is weak. FCM predicts *superlight* polarons and bipolarons in cuprates with a remarkably high superconducting critical temperature. Cuprate bipolarons are relatively light because they are $`intersite`$ rather than $`onsite`$ pairs due to the strong on-site repulsion, and because mainly $`c`$-axis polarized optical phonons are responsible for the in-plane mass renormalization. The relatively small mass renormalization of polaronic and bipolaronic carries in FCM has been confirmed numerically using the exact QMC Korn2 , cluster diagonalization feh3 and variational bon2 simulations.
(Bi)polarons describe many properties of cuprates alebook , in particular normal-state transport properties (section 2), the Nernst effect (section 3), and the normal state diamagnetism (section 4). The strong-coupling theory also provides an explanation for the phase separation and coexistence and describes the shape of resistive and magnetic transitions in manganites (section 5).
## II Normal state in-plane resistivity, Hall effect and magnetic susceptibility of cuprates in the bipolaron model
The low-energy FCM electronic structure of cuprates is shown in Fig.1 MOTT . Polaronic p-holes are bound into lattice inter-site singlets (A) or into singlets and triplets (B) (if spins are included in Eq.(2)) at any temperature. Above T<sub>c</sub> a charged bipolaronic Bose-liquid is non-degenerate and below $`T_c`$ phase coherence (ODLRO) of the preformed bosons sets in. The state above $`T_c`$ is perfectly ”normal” in the sense that the off-diagonal order parameter (i.e. the Bogoliubov-Gor’kov anomalous average $`(𝐫,𝐫^{})=\psi _{}\left(𝐫\right)\psi _{}(𝐫^{}`$) is zero above the resistive transition temperature $`T_c`$. Here $`\psi _,\left(𝐫\right)`$ annihilates electrons with spin $`,`$ at point $`𝐫`$. Triplet and singlet states are separated by the exchange energy $`J`$ which explains the spin gap observed in a number of NMR and neutron scattering experiments. There are also thermally excited single polarons in the model. Their density becomes comparable with the bipolaron density at the temperature $`T^{}`$ which is about half of the bipolaron binding energy $`\mathrm{\Delta }`$, in accordance with the experimentally observed crossover regime at $`T^{}>T_c`$ and the normal state pseudogaps in cuprates.
A nonlinear temperature dependence of the $`in`$-plane resistivity below $`T^{}`$, a temperature-dependent paramagnetic susceptibility, and a peculiar maximum in the Hall ratio well above $`T_c`$ have remained long-standing problems of cuprate physics. The bipolaron model provides their quantitative description alezavdzu . Thermally excited phonons and (bi)polarons are well decoupled in the strong-coupling regime of the electron-phonon interaction alebook , so the conventional Boltzmann kinetics for mobile polaronic and bipolaronic carries is applied. Here we use a ‘minimum’ bipolaron model Fig.1A, which includes the singlet bipolaron band and the spin 1/2 polaron band separated by $`T^{}`$, and the $`\tau `$approximation in weak electric $`𝐄`$ and magnetic fields, $`𝐁𝐄`$.
Bipolaron and single-polaron non-equilibrium distributions are found as
$$f\left(𝐤\right)=f_0\left(E\right)+\tau \frac{f_0}{E}𝐯\left\{𝐅+\mathrm{\Theta }𝐧\times 𝐅\right\},$$
(3)
where $`𝐯=E/𝐤,`$ $`𝐅=\stackrel{}{}\left(\mu 2e\varphi \right)`$, $`f_0\left(E\right)=\left[y^1\mathrm{exp}\left(E/T\right)1\right]^1`$ and the Hall angle $`\mathrm{\Theta }=\mathrm{\Theta }_b=2eB\tau _b/m_b`$ for bipolarons with the energy $`E=k^2/\left(2m_b\right)`$, and $`𝐅=\stackrel{}{}\left(\mu /2e\varphi \right)`$, $`f_0\left(E\right)=\left\{y^{1/2}\mathrm{exp}\left[\left(E+T^{}\right)/T\right]+1\right\}^1`$, $`E=k^2/\left(2m_p\right)`$, and $`\mathrm{\Theta }=\mathrm{\Theta }_p=eB\tau _p/m_p`$ for thermally excited polarons. Here $`m_b`$ and $`m_p`$ are the bipolaron and polaron mass, respectively, $`y=\mathrm{exp}\left(\mu /T\right),`$ $`\mu `$ is the chemical potential, and $`𝐧=𝐁/B`$ is a unit vector in the direction of the magnetic field. Eq.(3) is used to calculate the electrical resistivity and the Hall ratio as
$`\rho `$ $`=`$ $`{\displaystyle \frac{m_b}{4e^2\tau _bn_b\left(1+An_p/n_b\right)}},`$ (4)
$`R_H`$ $`=`$ $`{\displaystyle \frac{1+2A^2n_p/n_b}{2en_b\left(1+An_p/n_b\right)^2}},`$ (5)
where $`A=\tau _pm_b/\left(4\tau _bm_p\right)`$. The atomic densities of quasi two-dimensional carriers are found as
$`n_b={\displaystyle \frac{m_bT}{2\pi }}\left|\mathrm{ln}\left(1y\right)\right|,`$ (6)
$`n_p={\displaystyle \frac{m_pT}{\pi }}\mathrm{ln}\left[1+y^{1/2}\mathrm{exp}\left(T^{}/T\right)\right].`$ (7)
and the chemical potential is determined by doping $`x`$ using $`2n_b+n_p=xn_L`$, where $`n_L`$ is the number of carriers localised by disorder (here we take the lattice constant $`a=1`$).
Polarons are not degenerate. Their number remains small compared with twice the number of bipolarons, $`n_p/\left(2n_b\right)<0.2`$, in the relevant temperature range $`T<T^{}`$, so that
$$y1\mathrm{exp}\left(T_0/T\right),$$
(8)
where $`T_0=\pi \left(xn_L\right)/m_bT_c`$ is about the superconducting critical temperature of the (quasi)two-dimensional Bose gas. Because of this reason, the experimental $`T_c`$ was taken as $`T_0`$ in our fits. Using Eqs.(7,6,5) we obtain
$$R_H\left(T\right)=R_{H0}\frac{1+2A^2y^{1/2}\left(T/T_c\right)\mathrm{exp}\left(T^{}/T\right)}{\left[1+A\left(T/T_c\right)y^{1/2}\mathrm{exp}\left(T^{}/T\right)\right]^2},$$
(9)
where $`R_{H0}=\left[e\left(xn_L\right)\right]^1`$. If we assume that the number of localised carriers depends only weakly on temperature in underdoped cuprates since their average ionisation energy is sufficiently large, then $`R_{H0}`$ is temperature independent at $`T<T^{}`$. As proposed in Ref.BRAT the scattering rate at relatively high temperatures is due to inelastic collisions of itinerant carriers with those localised by disorder, so it is proportional to $`T^2`$. We also have to take into account the residual scattering of polarons off optical phonons, so that $`\tau ^1=aT^2+b\mathrm{exp}\left(\omega /T\right)`$, if the temperature is low compared with the characteristic phonon energy $`\omega `$. The relaxation times of each type of carriers scales with their charge $`e^{}`$ and mass as $`\tau _{p,b}m_{p,b}^{3/2}\left(e^{}\right)^2`$, so we estimate $`A=\left(m_b/m_p\right)^{5/2}6`$ if we take $`m_b2m_p`$ . As a result the in-plane resistivity is given by
$$\rho \left(T\right)=\rho _0\frac{\left(T/T_1\right)^2+\mathrm{exp}\left(\omega /T\right)}{\left[1+A\left(T/T_c\right)y^{1/2}\mathrm{exp}\left(T^{}/T\right)\right]},$$
(10)
where $`\rho _0=bm_b/\left[2e^2\left(xn_L\right)\right]`$ and $`T_1=\left(b/a\right)^{1/2}`$ are temperature independent. Finally, one can easily obtain the uniform magnetic susceptibility due to nondegenerate spin 1/2 polarons as AKM
$$\chi \left(T\right)=By^{1/2}\mathrm{exp}\left(T^{}/T\right)+\chi _0,$$
(11)
where $`B=\left(\mu _B^2m_p/\pi \right)`$, and $`\chi _0`$ is the magnetic susceptibility of the parent Mott insulator.
| $`\delta `$ | $`T_c`$ | $`\rho _0`$ | $`R_{H0}`$ | $`10^4B`$ | $`10^4\chi _0`$ | $`T^{}`$ | $`\omega `$ | $`T_1`$ |
| --- | --- | --- | --- | --- | --- | --- | --- | --- |
| | K | $`m\mathrm{\Omega }cm`$ | $`\frac{10^9m^3}{C}`$ | $`\frac{emu}{mole}`$ | $`\frac{emu}{mole}`$ | K | K | K |
| 0.05 | 90.7 | 1.8 | 0.45 | | | 144 | 447 | 332 |
| 0.12 | 93.7 | | | 2.6 | 2.1 | 155 | | |
| 0.19 | 87 | 3.4 | 0.63 | 4.5 | 1.6 | 180 | 477 | 454 |
| 0.23 | 80.6 | 5.7 | 0.74 | | | 210 | 525 | 586 |
| 0.26 | 78 | | | 5.4 | 1.5 | 259 | | |
| 0.28 | 68.6 | 8.9 | 0.81 | | | 259 | 594 | 786 |
| 0.38 | 61.9 | | | 7.2 | 1.4 | 348 | | |
| 0.39 | 58.1 | 17.8 | 0.96 | | | 344 | 747 | 1088 |
| 0.51 | 55 | | | 9.1 | 1.3 | 494 | | |
The present model numerically fits the Hall ratio, $`R_H\left(T\right)`$, the in-plane resistivity, $`\rho \left(T\right)`$, and the magnetic susceptibility $`\chi \left(T\right)`$ of $`YBa_2Cu_3O_{7\delta }`$ within the physically relevant range of all parameters (see Fig. 2 and the Table). The ratio of polaron and bipolaron mobilities $`A=7`$ used in all fits is close to the above estimate, and $`\chi _01.5\times 10^4emu/mole`$ is very close to the susceptibility of a slightly doped insulator coop . The maximum of $`R_H\left(T\right)`$ is due to the contribution of thermally excited polarons into transport, and the temperature dependence of the in-plane resistivity below $`T^{}`$ is due to this contribution and the combination of the carrier-carrier and carrier-phonon scattering. The characteristic phonon frequency from the resistivity fit (Table) decreases with doping and the pseudogap $`T^{}`$ shows the doping behaviour as observed in other independent experiments.
Notwithstanding our explanation of the Hall ratio, the in-plane resistivity and the bulk magnetic susceptibility might be not so convincing as a direct measurement of the double charge $`2e`$ on carriers in the normal state. In 1993, we discussed the thermal conductivity of preformed bosons NEV . The contribution from carriers to the thermal transport provided by the Wiedemann-Franz law depends strongly on the elementary charge as $`\left(e^{}\right)^2`$ and should be significantly suppressed if $`e^{}=2e`$. The Lorenz number, $`L`$, has been directly measured in $`YBa_2Cu_3O_{6.95}`$ by Zhang et al. zha using the thermal Hall conductivity. Remarkably, the measured value of $`L`$ just above $`T_c`$ was found just the same as predicted by the bipolaron model NEV , $`L0.15L_e`$, where $`L_e`$ is the conventional Fermi-liquid Lorenz number. The breakdown of the Wiedemann-Franz law has been also explained in the framework of the bipolaron model leeale .
## III Normal-state Nernst effect
In disagreement with the weak-coupling BCS and the strong-coupling bipolaron theories a significant fraction of research in the field of high-temperature superconductivity suggests that the superconducting transition is only a phase ordering while the superconducting order parameter $`(𝐫,𝐫^{})`$ remains nonzero above the resistive $`T_c`$. One of the key experiments supporting this viewpoint is the large Nernst signal observed in the normal (i.e. resistive) state of cuprates (see Ref. xu ; cap ; cap2 and references therein). Some authors xu ; ong claim that numerous resistive determinations of the upper critical field, $`H_{c2}\left(T\right)`$ in cuprates have been misleading since the Nernst signal xu and the diamagnetic magnetization ong imply that $`H_{c2}\left(T\right)`$ remains large at $`T_c`$ and above. They propose a ”vortex scenario”, where the long-range phase coherence is destroyed by mobile vortices, but the amplitude of the off-diagonal order parameter remains finite and the Cooper pairing with a large binding energy exists well above $`T_c`$ supporting the so-called ”preformed Cooper-pair” or ”phase fluctuation” model kiv . The model is based on the assumption that the superfluid density is small compared with the normal carrier density in cuprates. These interpretations seriously undermine many theoretical and experimental works on superconducting cuprates, which consider the state above $`T_c`$ as perfectly normal with no off-diagonal order, either long or short.
We believe that the vortex (or phase fluctuation) scenario contradicts straightforward resistive and other measurements, and it is theoretically inconsistent. This scenario is impossible to reconcile with the extremely sharp resistive transitions at $`T_c`$ in high-quality underdoped, optimally doped and overdoped cuprates. For example, the in-plane and out-of-plane resistivity of $`Bi2212`$, where the anomalous Nernst signal has been measured xu , is perfectly ”normal” above $`T_c`$, Fig.3, showing only a few percent positive or negative magnetoresistance zavale .
Both in-plane buc ; mac0 ; boz ; fra ; gan and out-of-plane alezavnev ; out ; out2 resistive transitions of high-quality samples are sharp and remain sharp in the magnetic field providing a reliable determination of the genuine $`H_{c2}\left(T\right)`$. The vortex entropy cap estimated from the Nernst signal is an order of magnitude smaller than the difference between the entropy of the superconducting state and the extrapolated entropy of the normal state obtained from the specific heat. The preformed Cooper-pair model kiv is incompatible with a great number of thermodynamic, magnetic, and kinetic measurements, which show that only holes (density x), doped into a parent insulator are carriers *both* in the normal and the superconducting states of cuprates. The assumption kiv that the superfluid density is small compared with the normal-state carrier density is also inconsistent with the theorem leg , which proves that the number of supercarriers at $`T=0`$K should be the same as the number of normal-state carriers in any clean superfluid.
Recently we described the unusual Nernst signal in cuprates in a different manner as the normal state phenomenon alezav . Here we extend our description to cuprates with very low doping level accounting for their Nernst signal, the thermopower and the insulating-like in-plane low temperature resistance xu ; cap ; cap2 .
Thermomagnetic effects appear in conductors subjected to a longitudinal temperature gradient $`_xT`$ in $`x`$ direction and a perpendicular magnetic field in $`z`$ direction. The transverse Nernst-Ettingshausen effect nernst (here the Nernst effect) is the appearance of a transverse electric field $`E_y`$ in the third direction. When bipolarons are formed in the strong-coupling regime, the chemical potential is negative, Eq.(8). It is found in the impurity band just below the mobility edge at $`T>T_c`$. Carriers, localised below the mobility edge contribute to the longitudinal transport together with the itinerant carriers in extended states above the mobility edge. Importantly the contribution of localised carriers of any statistics to the *transverse* transport is normally small ell since a microscopic Hall voltage will only develop at junctions in the intersections of the percolation paths, and it is expected that these are few for the case of hopping conduction among disorder-localised states mott2 . Even if this contribution is not negligible, it adds to the contribution of itinerant carriers to produce a large Nernst signal, $`e_y(T,B)E_y/_xT`$, while it reduces the thermopower $`S`$ and the Hall angle $`\mathrm{\Theta }`$. This unusual ”symmetry breaking” is completely at variance with ordinary metals where the familiar ”Sondheimer” cancelation sond makes $`e_y`$ much smaller than $`S\mathrm{tan}\mathrm{\Theta }`$ because of the electron-hole symmetry near the Fermi level. Such behaviour originates in the ”sign” (or ”$`pn`$”) anomaly of the Hall conductivity of localised carriers. The sign of their Hall effect is often $`opposite`$ to that of the thermopower as observed in many amorphous semiconductors ell and described theoretically fri .
The Nernst signal is expressed in terms of the kinetic coefficients $`\sigma _{ij}`$ and $`\alpha _{ij}`$ as
$$e_y=\frac{\sigma _{xx}\alpha _{yx}\sigma _{yx}\alpha _{xx}}{\sigma _{xx}^2+\sigma _{xy}^2},$$
(12)
where the current density is given by $`j_i=\sigma _{ij}E_j+\alpha _{ij}_jT`$. When the chemical potential $`\mu `$ is at the mobility edge, the localised carriers contribute to the transport, so $`\sigma _{ij}`$ and $`\alpha _{ij}`$ in Eq.(12) can be expressed as $`\sigma _{ij}^{ext}+\sigma _{ij}^l`$ and $`\alpha _{ij}^{ext}+\alpha ^lij`$, respectively. Since the Hall mobility of carriers localised below $`\mu `$, $`\sigma _{yx}^l`$, has the sign opposite to that of carries in the extended states above $`\mu `$, $`\sigma _{yx}^{ext}`$, the sign of the off-diagonal Peltier conductivity $`\alpha _{yx}^l`$ should be the same as the sign of $`\alpha _{yx}^{ext}`$. Then neglecting the magneto-orbital effects in the resistivity (since $`\mathrm{\Theta }1`$ xu ) we obtain
$$S\mathrm{tan}\mathrm{\Theta }\frac{\sigma _{yx}\alpha _{xx}}{\sigma _{xx}^2+\sigma _{xy}^2}\rho \left(\alpha _{xx}^{ext}\left|\alpha _{xx}^l\right|\right)\left(\mathrm{\Theta }^{ext}\left|\mathrm{\Theta }^l\right|\right)$$
(13)
and
$$e_y\rho \left(\alpha _{yx}^{ext}+\left|\alpha _{yx}^l\right|\right)S\mathrm{tan}\mathrm{\Theta },$$
(14)
where $`\mathrm{\Theta }^{ext}\sigma _{yx}^{ext}/\sigma _{xx}`$, $`\mathrm{\Theta }^l\sigma _{yx}^l/\sigma _{xx}`$, and $`\rho =1/\sigma _{xx}`$ is the resistivity.
Clearly the model, Eqs.(13,14) can account for a low value of $`S\mathrm{tan}\mathrm{\Theta }`$ compared with a large value of $`e_y`$ in some underdoped cuprates xu ; cap2 due to the sign anomaly. Even in the case when localised bosons contribute little to the conductivity their contribution to the thermopower $`S=\rho (\alpha _{xx}^{ext}|\alpha _{xx}^l\left|\right))`$ could almost cancel the opposite sign contribution of itinerant carriers alezav . Indeed the longitudinal conductivity of itinerant two-dimensional bosons, $`\sigma ^{ext}_0𝑑EE𝑑f\left(E\right)/𝑑E`$ diverges logarithmically when $`\mu `$ in the Bose-Einstein distribution function $`f\left(E\right)=\left[\mathrm{exp}\left(\left(E\mu \right)/T\right)1\right]^1`$ goes to zero and the relaxation time $`\tau `$ is a constant. At the same time $`\alpha _{xx}^{ext}_0𝑑EE\left(E\mu \right)𝑑f\left(E\right)/𝑑E`$ remains finite, and it could have the magnitude comparable with $`\alpha _{xx}^l`$. Statistics of bipolarons gradually changes from Bose to Fermi statistics with lowering energy across the mobility edge because of the Coulomb repulsion of bosons in localised states alegile . Hence one can use the same expansion near the mobility edge as in ordinary amorphous semiconductors to obtain the familiar textbook result $`S=S_0T`$ with a constant $`S_0`$ at low temperatures mott3 . The model becomes particularly simple, if we neglect the localised carrier contribution to $`\rho `$, $`\mathrm{\Theta }`$ and $`\alpha _{xy}`$, and take into account that $`\alpha _{xy}^{ext}B/\rho ^2`$ and $`\mathrm{\Theta }^{ext}B/\rho `$ in accordance with the Boltzmann theory. Then Eqs.(13,14) yield
$$S\mathrm{tan}\mathrm{\Theta }T/\rho $$
(15)
and
$$e_y(T,B)\left(1T/T_1\right)/\rho .$$
(16)
According to our earlier suggestion alelog the insulating-like low-temperature dependence of $`\rho \left(T\right)`$ in underdoped cuprates originates from the elastic scattering of nondegenerate itinerant carriers off charged impurities. As in section 2 we assume here that the carrier density is temperature independent at low temperatures in agreement with the temperature-independent Hall effect per . The relaxation time of nondegenerate carriers depends on temperature as $`\tau T^{1/2}`$ for scattering off short-range deep potential wells, and as $`T^{1/2}`$ for very shallow wells alelog . Combining both scattering rates we obtain
$$\rho =\rho _0\left[\left(T/T_2\right)^{1/2}+\left(T_2/T\right)^{1/2}\right].$$
(17)
Eq.(17) with $`\rho _0=0.236`$ m$`\mathrm{\Omega }`$cm and $`T_2=44.6`$K fits extremely well the experimental insulating-like normal state resistivity of underdoped La<sub>1.94</sub> Sr<sub>0.06</sub>CuO<sub>4</sub> in the whole low-temperature range from 2K up to 50K, Fig.4, as revealed in the field $`B=12`$ Tesla cap ; cap2 . Another high quality fit can be obtained combining the Brooks-Herring formula for the 3D scattering off charged impurities, as proposed in Ref.kast for almost undoped $`LSCO`$, or the Coulomb scattering in 2D ($`\tau T`$) and a temperature independent scattering rate off neutral impurities with the carrier exchange erg similar to the scattering of slow electrons by hydrogen atoms. Importantly our expressions (15,16) for $`S\mathrm{tan}\mathrm{\Theta }`$ and $`e_y`$ do not depend on the particular scattering mechanism. Taking into account the excellent fit of Eq.(17) to the experiment, they can be parameterized as
$$S\mathrm{tan}\mathrm{\Theta }=e_0\frac{\left(T/T_2\right)^{3/2}}{1+T/T_2},$$
(18)
and
$$e_y(T,B)=e_0\frac{\left(T_1T\right)\left(T/T_2\right)^{1/2}}{T_2+T},$$
(19)
where $`T_1`$ and $`e_0`$ are temperature independent.
In spite of all simplifications, the model describes remarkably well both $`S\mathrm{tan}\mathrm{\Theta }`$ and $`e_y`$ measured in La<sub>1.94</sub> Sr<sub>0.06</sub>CuO<sub>4</sub> with a $`single`$ fitting parameter, $`T_1=50`$K using the experimental $`\rho \left(T\right)`$. The constant $`e_0=2.95`$ $`\mu `$V/K scales the magnitudes of $`S\mathrm{tan}\mathrm{\Theta }`$ and $`e_y`$. The magnetic field $`B=12`$ Tesla destroys the superconducting state of the low-doped La<sub>1.94</sub> Sr<sub>0.06</sub>CuO<sub>4</sub> down to $`2`$K, Fig.4, so any residual superconducting order above $`2`$K is clearly ruled out, while the Nernst signal, Fig.5, is remarkably large. The coexistence of the large Nernst signal and a nonmetallic resistivity is in sharp disagreement with the vortex scenario, but in agreement with our model. Taking into account the field dependence of the conductivity of localised carriers, the phonon-drug effect, and their contribution to the transverse magnetotransport can well describe the magnetic field dependence of the Nernst signal alezav and improve the fit in Fig.5 at the expense of the increasing number of fitting parameters.
## IV Normal state diamagnetism in cuprates
A number of experiments (see, for example, mac ; jun ; hof ; nau ; igu ; ong and references therein), including torque magnetometries, showed enhanced diamagnetism above $`T_c`$, which has been explained as the fluctuation diamagnetism in quasi-2D superconducting cuprates (see, for example Ref. hof ). The data taken at relatively low magnetic fields (typically below 5 Tesla) revealed a crossing point in the magnetization $`M(T,B)`$ of most anisotropic cuprates (e.g. $`Bi2212`$), or in $`M(T,B)/B^{1/2}`$ of less anisotropic $`YBCO`$ jun . The dependence of magnetization (or $`M/B^{1/2}`$) on the magnetic field has been shown to vanish at some characteristic temperature below $`T_c`$. However the data taken in high magnetic fields (up to 30 Tesla) have shown that the crossing point, anticipated for low-dimensional superconductors and associated with superconducting fluctuations, does not explicitly exist in magnetic fields above 5 Tesla nau .
Most surprisingly the torque magnetometery mac ; nau uncovered a diamagnetic signal somewhat above $`T_c`$ which increases in magnitude with applied magnetic field. It has been linked with the Nernst signal and mobile vortexes in the normal state of cuprates ong . However, apart from the inconsistences mentioned above, the vortex scenario of the normal-state diamagnetism is internally inconsistent. Accepting the vortex scenario and fitting the magnetization data in $`Bi2212`$ with the conventional logarithmic field dependence ong , one obtains surprisingly high upper critical fields $`H_{c2}>120`$ Tesla and a very large Ginzburg-Landau parameter, $`\kappa =\lambda /\xi >450`$ even at temperatures close to $`T_c`$. The in-plane low-temperature magnetic field penetration depth is $`\lambda =200`$ nm in optimally doped $`Bi2212`$ (see, for example tal ). Hence the zero temperature coherence length $`\xi `$ turns out to be about the lattice constant, $`\xi =0.45`$nm, or even smaller. Such a small coherence length rules out the ”preformed Cooper pairs” kiv , since the pairs are virtually not overlapped at any size of the Fermi surface in $`Bi2212`$ . Moreover the magnetic field dependence of $`M(T,B)`$ at and above $`T_c`$ is entirely inconsistent with what one expects from a vortex liquid. While $`M\left(B\right)`$ decreases logarithmically at temperatures well below $`T_c`$, the experimental curves mac ; nau ; ong clearly show that $`M\left(B\right)`$ increases with the field at and above $`T_c`$ , just opposite to what one could expect in the vortex liquid. This significant departure from the London liquid behavior clearly indicates that the vortex liquid does not appear above the resistive phase transition mac .
Some time ago we explained the anomalous diamagnetism in cuprates as the Landau normal-state diamagnetism of preformed bosons den . The same model predicted the unusual upper critical field aleH observed in many superconducting cuprates buc ; mac0 ; boz ; fra ; gan ; alezavnev ; ZAV . Here we extend the model to high magnetic fields taking into account the magnetic pair-breaking of singlet bipolarons and the anisotropy of the energy spectrum.
When the strong magnetic field is applied perpendicular to the copper-oxygen plains the quasi-2D bipolaron energy spectrum is quantized as
$$E_\alpha =\omega \left(n+1/2\right)+2t_c\left[1\mathrm{cos}\left(k_zd\right)\right],$$
(20)
where $`\omega =2eB/m_b`$, $`n=0,1,2,\mathrm{}`$, and $`t_c`$, $`k_z`$, $`d`$ are the hopping integral, the momentum and the lattice period perpendicular to the planes. Quantum numbers $`\alpha `$ also include the momentum along one of the in-plane directions. Expanding the Bose-Einstein distribution function in powers of $`exp\left[\left(\mu E_\alpha \right)/T\right]`$ with the negative $`\mu `$ one can readily obtain (after summation over $`n`$) the boson density
$$n_b=\frac{eB}{\pi d}\underset{k=1}{\overset{\mathrm{}}{}}I_0\left(2t_ck/T\right)\frac{\mathrm{exp}\left[\left(\stackrel{~}{\mu }2t_c\right)k/T\right]}{1\mathrm{exp}\left(\omega k/T\right)},$$
(21)
and the magnetization
$`M(T,B)`$ $`=`$ $`n_b\mu _b+{\displaystyle \frac{eT}{\pi d}}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}I_0\left(2t_ck/T\right){\displaystyle \frac{\mathrm{exp}\left[\left(\stackrel{~}{\mu }2t_c\right)k/T\right]}{1\mathrm{exp}\left(\omega k/T\right)}}`$ (22)
$`\times `$ $`\left({\displaystyle \frac{1}{k}}{\displaystyle \frac{\omega \mathrm{exp}\left(\omega k/T\right)}{T\left[1\mathrm{exp}\left(\omega k/T\right)\right]}}\right),`$
where $`\mu _b=e/m_b`$, $`\stackrel{~}{\mu }=\mu \omega /2`$ and $`I_0\left(x\right)`$ is the modified Bessel function. At low temperatures $`T0`$ Schafroth’s result sha is recovered, $`M(0,B)=n_b\mu _b`$. The magnetization of charged bosons is field-independent at low temperatures. At high temperatures, $`TT_c`$ the chemical potential has a large magnitude , so we can keep only terms with $`k=1`$ in Eqs.(21,22) to obtain
$$M(T,B)=n_b\mu _b+\frac{Tn_b}{B}\left(1\frac{\omega \mathrm{exp}\left(\omega /T\right)}{T\left[1\mathrm{exp}\left(\omega /T\right)\right]}\right).$$
(23)
The experimental conditions are such that $`T\omega `$ when $`T`$ is of the order of $`T_c`$ or higher, so that
$$M(T,B)=n_b\mu _b\frac{\omega }{6T},$$
(24)
which is the Landau orbital diamagnetism of nondegenerate carriers. The bipolaron in-plane mass in cuprates is about $`m_b10m_e`$ alebook . Using this mass yields $`M(0,B)2000`$ A/m with the bipolaron density $`n_b=10^{21}`$ cm<sup>-3</sup>. Then the magnitude and the field/temperature dependence of $`M(T,B)`$ near and above $`T_c`$ are about the same as experimentally observed in Refs nau ; ong . The pseudogap temperature $`T^{}`$ depends on the magnetic field predominantly because of the magnetic-field splitting of the single-polaron band in Fig.1. As a result the bipolaron density depends on the field (as well as on temperature) near $`T_c`$ as
$$n_b(T,B)=n_b(T_c,0)\left[1+\left(T_cT\right)/\stackrel{~}{T}_0\left(B/B_0\right)^\beta \right],$$
(25)
where $`\stackrel{~}{T}_0`$ and $`B_0`$ are constants depending on $`T^{}`$, $`\beta =2`$ if the polaron spectrum is spin-degenerate, and $`\beta =1`$ if the spin degeneracy is removed by the crystal field already in the absence of the external field.
Theoretical temperature and field dependencies of $`M(T,B)`$, Eq.(22) agree qualitatively with the experimental curves in $`Bi2212`$ nau ; ong , if the depletion of the bipolaron density, Eq.(25) is taken into account. The depletion of $`n_b`$ accounts for the absence of the crossing point in $`M(T,B)`$ at high magnetic fields. Nevertheless a quantitative fit to experimental $`M(T,B)`$ curves using $`\stackrel{~}{T}_0`$ and $`B_0`$ as the fitting parameters is premature. The experimental diamagnetic magnetization has been extracted from the total magnetization assuming that the normal state paramagnetic contribution remains temperature-independent at all temperatures nau ; ong . This assumption is inconsistent with a great number of NMR and the Knight shift measurements, and even with the preformed Cooper-pair model itself. The Pauli spin-susceptibility has been found temperature-dependent in these experiments revealing a normal-state pseudogap, contrary to the assumption. Hence the experimental diamagnetic $`M(T,B)`$ nau ; ong has to be corrected by taking into account the temperature dependence of the spin paramagnetism at relatively low temperatures.
## V Phase coexistence and resistivity near the ferromagnetic transition in manganites
Ferromagnetic oxides, in particular manganese perovskites, show a huge magnetoresistance near the ferromagnetic transition. The resistivity change is so large that it could not compare with any other forms of magnetoresistance. The effect observed in these materials was therefore named ’colossal’ magnetoresistance (CMR) to distinguish it from the giant magnetoresistance observed in magnetic multilayers. The discovery raised expectations of a new generation of magnetic devices, and launched a frenetic scientific race to understand the cause of the effect. Significant progress has been made in understanding their properties, but new questions have arisen. The ferromagnetic metal-insulator transition in manganites has long been thought as the consequence of the so-called double exchange mechanism (DEX), which results in a varying bandwidth of electrons in the Mn<sup>3+</sup> d-shell as a function of temperature dex . More recently it has been noticed mil that the effective spin-exchange interaction of the double-exchange model cannot account for CMR alone. In fact there is strong experimental evidence for exceptionally strong e-ph interactions in doped manganites from the optical data (see section 2), the giant isotope effect zhao1 , the Arrhenius behaviour of the drift and Hall mobilities emi0 in the paramagnetic phase above the Curie temperature, $`T_m`$, etc. Therefore Ref. mil and some subsequent theoretical studies combined DEX with the Jahn-Teller e-ph interaction in d-orbitals arriving at the conclusion that the low-temperature ferromagnetic phase is a spin-polarised metal, while the paramagnetic phase is a polaronic insulator.
However, some low-temperature optical opt , electron-energy-loss (EELS) eels , photoemissionarpes and thermoelectric thermo measurements showed that the ferromagnetic phase of manganites is not a conventional metal. In particular, broad incoherent spectral features and a pseudo-gap in the excitation spectrum were observed. EELS confirmed that manganites were charge-transfer doped insulators having p-holes as current carriers rather than d Mn<sup>3+</sup> electrons. Photoemission and x-ray absorption spectroscopies of La<sub>1-x</sub>Sr<sub>x</sub>MnO<sub>3</sub> also showed that the itinerant holes doped into LaMnO<sub>3</sub> are of oxygen p-character. CMR has been observed in the ferromagnetic pyrochlore manganite Tl<sub>2</sub>Mn<sub>2</sub>O<sub>7</sub> ram , which has neither the mixed valence for DEX magnetic interaction nor the Jahn-Teller cations such as Mn<sup>3+</sup>.
These and other observations zhao2 , in particular the fact that some samples of ferromagnetic manganites manifest an insulating-like optical conductivity at all temperatures emi , clearly rule out DEX as the mechanism of CMR. They led us to a novel theory of ferromagnetic/paramagnetic phase transition and CMR based on the so-called current-carrier density collapse (CCDC) alebra2 . In CCDC p-holes are bound into heavy bipolarons above $`T_m`$ due to the Fröhlich e-ph interaction, Eq.(2). The resistivity peak and CMR are the result of the magnetic pair-breaking below $`T_m`$, Fig.6, caused by the $`pd`$ spin-exchange interaction, $`J_{pd}`$, which described as
$$H_{pd}=\left(2N\right)^1\underset{𝐧,𝐦}{}J_{pd}\widehat{S}_𝐦^z\left(c_𝐧^{}c_𝐧c_𝐧^{}c_𝐧\right).$$
(26)
Here $`\widehat{S}_𝐦^z`$ is the z-component of $`Mn^{3+}`$ spin on site $`𝐦`$, and $`N`$ is the total number of sites.
Different from cuprates hole bipolarons are much heavier in manganites because the e-ph Fröhlich interaction is stronger and the band structure is less anisotropic. They are readily localised by disorder, so only thermally excited single extended polarons conduct in the paramagnetic phase. With temperature lowering single polarons polarize manganese spins at $`T_m`$ via $`J_{pd}`$, and the spin polarization of manganese ions breaks the bipolaronic singlets creating a spin-polarized $`polaronic`$ conductor. CCDC explained the resistivity peak and CMR in the experimental range of external magnetic fields alebra2 ; tai . More recently, the theory has been further confirmed experimentally. In particular, the oxygen isotope effect has been observed in the low-temperature resistivity of La<sub>0.75</sub>Ca<sub>0.25</sub>MnO<sub>3</sub> and Nd<sub>0.7</sub>Sr<sub>0.3</sub>MnO<sub>3</sub> and explained by CCDC with polaronic carriers in the ferromagnetic phase alezhao . The current-carrier density collapse has been directly observed using the Hall data in La<sub>0.67</sub>Ca<sub>0.33</sub>MnO<sub>3</sub> and La<sub>0.67</sub>Sr<sub>0.33</sub>MnO<sub>3</sub> hall . And the first order phase transition at $`T_m`$, predicted by the theory alebra2 , has been firmly established in the specific heat measurements phil . On the other hand, the resistivity and the magnetization of some samples of La<sub>0.7</sub>Ca<sub>0.3</sub>Mn<sub>1-x</sub>Ti<sub>x</sub>O<sub>3</sub> showed a more gradual (second-order like) transition china . Also the coexistence of ferromagnetic and paramagnetic phases near the Curie temperature observed in tunneling tun and other experiments has not yet been addressed in the framework of CCDC. Here we argue that the diagonal disorder, which is inevitable with doping, explains both the phase coexistence and the resistivity/magnetization shape near the transition.
The mean-field equations alebra2 describing the single polaron density $`n`$, p-hole polaron $`m`$ and manganese $`\sigma `$ magnetizations, and the chemical potential $`\mu =T\mathrm{ln}y`$ can be easily generalized by taking into account the random distribution of the bipolaron binding energy $`\delta =\mathrm{\Delta }/\left(2J_{pd}\right)`$ across the sample as
$`n_i=6y\mathrm{cosh}\left(\sigma _i/t\right),`$ (27)
$`m_i=n_i\mathrm{tanh}\left(\sigma _i/t\right),`$ (28)
$`\sigma _i=B_2\left(m_i/2t\right),`$ (29)
$`y^2={\displaystyle \frac{xn_i}{18}}\mathrm{exp}\left(2\delta _i/t\right),`$ (30)
where $`t=T/J_{pd}`$ is the reduced temperature, $`B_S`$ is the Brillouin function, $`x`$ is the number of delocalised holes at zero temperature in p-orbital states, which are $`3`$-fold degenerate. The subscript $`i`$ means different parts of the sample with different $`\delta _i`$ because of disorder. While averaging these transparent equations over a random distribution of $`\delta _i`$ is rather cumbersome, one can apply a simplified approach using the fact that the phase transition in a homogeneous system is of the first order in a wide range of $`\delta `$ alebra2 . Taking $`\sigma _i=\mathrm{\Theta }(T_{mi}T)`$ and $`n_i=x\mathrm{\Theta }(T_{mi}T)+\sqrt{2x}\mathrm{exp}(\mathrm{\Delta }/\left(2T\right))\mathrm{\Theta }(TT_{mi})`$ and averaging both quantities with the Gaussian distribution of random $`T_{mi}`$s around the experimental $`T_m`$ we obtain the averaged manganese magnetization
$$\sigma \left(T\right)=\frac{1}{2}erfc\left(\frac{TT_m}{\mathrm{\Gamma }}\right)$$
(31)
and the resistivity, $`\rho 1/n`$, near the transition
$`1/\rho \left(T\right)`$ $``$ $`erfc\left({\displaystyle \frac{TT_m}{\mathrm{\Gamma }}}\right)`$ (32)
$`+`$ $`\left(2/x\right)^{1/2}e^{\mathrm{\Delta }/2T}erfc\left({\displaystyle \frac{T_mT}{\mathrm{\Gamma }}}\right).`$ (33)
Here $`\mathrm{\Delta }`$ is the average bipolaron binding energy, $`\mathrm{\Theta }\left(y\right)=1`$ for $`y>0`$ and zero for $`y<0`$ , and $`erfc\left(y\right)=\left(2/\pi ^{1/2}\right)_y^{\mathrm{}}𝑑y\mathrm{exp}\left(y^2\right)`$. CCDC with disorder, Eq.(29) fits nicely the experimental resistivity china near the transition with physically reasonable parameters $`\mathrm{\Gamma }=28`$K, $`\mathrm{\Delta }=1600`$K, $`T_m=102`$K, and $`x=0.1`$, Fig.7. A random distribution of transition temperatures with the width $`\mathrm{\Gamma }`$ across the sample caused by the randomness of the bipolaron binding energy is responsible for the phase coexistence near the transition tun .
In summary, the strong-coupling bipolaron extension of the BCS theory accounts for the kinetic properties of superconducting cuprates including the temperature-dependent spin susceptibility, the nonlinear in-plane resistivity, the maximum in the Hall effect, the normal-state Nernst signal and the diamagnetism near and above $`T_c`$. CMR and ferromagnetism of ferromagnetic oxides can be well explained by the current-carrier density collapse in the framework of the same theory including the exchange magnetic interaction of p-holes with the manganese spins and disorder effects.
I thank A.M. Bratkovsky, J.P. Hague, V.V. Kabanov, P.E. Kornilovitch, J.H. Samson, P.E. Spencer, and V.N. Zavaritsky for collaboration and valuable discussions. The work was supported by EPSRC (UK) (grant EP/C518365/1).
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# Early re-brightening of the afterglow of GRB 050525a Based on observations performed with TAROT at the Calern observatory
## 1 Introduction
GRB 050525a was a very bright gamma-ray burst (GRB) detected on the 25<sup>th</sup> of May 2005 at 00:02:53.3 UT (hereafter $`t_{\mathrm{trig}}`$) by the BAT instrument on the Swift spacecraft (trigger=130088, Band et al. band2005 (2005)). The gamma-ray light curve shows that GRB 050525a is a multipeak GRB, with an emission lasting approximately 10 sec above 50 keV. The fluence of GRB 050525a in the range 20-1000 keV is 7.84 10<sup>-5</sup> erg/cm<sup>2</sup>, and its peak energy is $`E_{\mathrm{peak}}`$ = 84 keV (Golenetskii et al. golenetskii2005 (2005)). Spectroscopic observations performed 11 hours after the GRB revealed absorption lines from the host galaxy at a redshift z=0.606 (Foley et al. foley2005 (2005)). At that redshift, and adopting a flat cosmology with $`\mathrm{\Omega }_m=0.3`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$, and h$`{}_{0}{}^{}=0.65`$, the isotropic-equivalent energy of GRB 050525a in the range 1 keV to 10 MeV is $`E_{\mathrm{iso}}=12.6`$ $`10^{52}`$ erg. The intrinsic peak energy of the time integrated spectrum is 135$`\pm `$2 keV.
In this letter we report the early optical observations of the GRB 050525a afterglow, performed with the robotic TAROT observatory. The Gamma-ray bursts Coordinates Network (GCN) notice, providing celestial coordinates to ground stations, was send at 00:08:48 (Band et al. band2005 (2005)), too late to detect the hypothetical optical prompt emission. The first image of TAROT started at 00:08:52.1 UT, 5min59s after the GRB. The afterglow was detected on all images taken until the end of the night at the TAROT observatory (02:19 UT) at the coordinate quoted by Rykoff et al.Rykoff2005 (2005): R.A. 18h 32m 32.76s and Dec. +26$`{}_{}{}^{}20_{}^{}22.65^{\prime \prime }`$ (J2000.0). These data provide a continuous follow-up from 6 to 136 minutes after the GRB. In this paper we show that the classical exponential decay was perturbated at $`tt_{\mathrm{trig}}`$ 33 min by a re-brightening event.
Section 2 describes the technical details of the TAROT observations and of data reductions. In section 3, we compare our early time observations of GRB 050525a with those of other bursts with early optical observations (t $``$ 0.01-0.1 day) and dense sampling. In section 4 we discuss the theoretical interpretations which have been proposed to explain the early re-brightening of GRB optical afterglows.
## 2 TAROT observations
TAROT is a fully autonomous 25 cm aperture telescope installed at the Calern observatory (Observatoire de la Cote d’Azur - France). This telescope is devoted to very early observations of GRB optical counterparts. A technical description of TAROT can be read in Bringer et al. (Bringer99 (1999)) and in Bringer et al. (bringer2001 (2001)). The CCD camera is a commercial Andor based on a Marconi 4240 chip and is placed at the newtonian focus. The spatial sampling is 3.3 arcsec/pixel. The field of view is $`1.86\mathrm{°}`$. The readout noise is 9 electrons rms. The readout time is 5 seconds (to read the entire 2048x2048 matrix with no binning).
The first image was taken less than 4 seconds after the position of GRB 050525a was provided by the GCN. A series of 76 images of various exposure times (15, 30, 60, 120 s) were performed without any filter (hereafter clear filter). On-line preprocessing software enable to provide calibrated images less than 2 minutes after they were taken (corrected by dark, flat and astrometrically calibrated from the USNO-A1.0 catalog).
On the first 31 images, the afterglow is bright enough to be measured with a good accuracy on individual images. Later images were co-added to increase the signal to noise. Due to the decreasing of flux during exposures, the mean date of an observation, $`T`$, is not the middle of exposure; it must be interpolated between $`t_1`$ (start of the first frame) and $`t_2`$ (end of the last frame) such that the flux $`f`$ verify:
$`{\displaystyle _{t_1}^{t_2}}f(t)dt=f(T){\displaystyle _{t_1}^{t_2}}dt`$ (1)
Considering afterglow decay flux law: $`f(t)t^\alpha `$, and assuming all times counted since $`t_{\mathrm{trig}}`$,
$`T={\displaystyle \frac{t_2t_1}{\mathrm{ln}(t_2/t_1)}}`$ if $`\alpha `$=+1 (2)
$`T=\left[{\displaystyle \frac{(t_2t_1)(1\alpha )}{t_2{}_{}{}^{1\alpha }t_1^{1\alpha }}}\right]^{1/\alpha }`$ if $`\alpha `$+1 (3)
$`\alpha `$ and $`T`$ values are computed by iterations. Initial $`T`$ values are computed taking $`\alpha `$=+1. Then the fit of the first light curve refines the $`\alpha `$ value. The second iteration is enough to converge (see Table 1 and Fig. 1).
The presence of the star USNO-B1 1163-0325216 (R$``$16.6), at about 15 arcsec north-west from the afterglow, perturbated the photometry, especially when the afterglow fades. The set of magnitudes provided by Klotz et al. (klotz2005 (2005)) was affected by this effect, which leads to a false plateau for dates 40 min after the trigger. To eliminate this effect, we computed an image (designated hereafter mask) of the field which does not show the afterglow. The mask is synthetized in two steps: first we oversampled the 76 images by a factor three and stack them to synthetize the sum. In this image, the afterglow and the star are well separated. The second step is to clear only the afterglow spot in order to synthetize the mask. The mask was normalized in flux for each image and was substracted, leading to images where only the afterglow appears (and some residuals of bright stars). Then we can extract the magnitude of the afterglow avoiding problems of contamination by the nearby star.
The GRB 030329 afterglow shows there is no colour effect during the first phases of decay (Zeh et al. 2003a , 2003b ). Taking this advantage, we performed differential photometry with two stars as reference: USNO-B1 1163-0325130 (R=11.25) and USNO-B1 1163-0325158 (R=14.00). First, we verified that the magnitude of reference stars does not vary. As we have no information about the afterglow colours (we used no filter), we do not obtain directly R magnitudes. We designated by CR, the unfiltered magnitudes calibrated by USNO-B1 R magnitudes of reference stars. As the airmass only varied from 1.13 to 1.05 during measurements, its effect on colours is assumed to be negligible (presumatly lower than 0.05 mag, much less than other uncertainties). As a consequence, it is possible to convert our CR magnitudes into standard R magnitudes by a simple offset (estimated lower than 0.2 magnitude, depending on the intrinsic colour differences between afterglow and reference stars). Anyway, the CR magnitudes allows the compute decay slope parameter $`\alpha `$.
Light curve (Fig. 1) shows two parts separated at $`tt_{\mathrm{trig}}`$ 33 min. Fits of slopes give $`\alpha _1=1.14\pm 0.07`$ and $`\alpha _2=1.23\pm 0.27`$. Uncertainties are such that slopes are not significantly differents. The most important remark is the offset of about 0.65 magnitude (nearly a factor two) between the two parts. We can join the two curves by two extreme paths: i) a sharp re-brightening of 0.65 magnitude in less than 3 minute centered at $`tt_{\mathrm{trig}}=`$ 32.8 min, ii) a plateau (flat re-brightening) of CR $``$ 16.7, begining at $`tt_{\mathrm{trig}}`$ 26 min and finishing at $`tt_{\mathrm{trig}}`$ 43 min. The magnitude uncertainties of our measurements are too large to discriminate one assumption to the other.
On Fig. 1, we reported a measurement ($``$) obtained by Malesani et al. (malesani2005 (2005)). They used a large aperture telescope and a R filter. They found the afterglow 0.5 magnitude brighter than CR TAROT value. Less than 0.2 mag. of this offset can be due to the CR-R differences as explained in section 2. The remaining 0.3 mag may be due to the substraction of the USNO-B1 1163-0325216 nearby star. Due to a much larger aperture of telescope, Malesani mesurement is probably best accurate than the two last measurements obtained with TAROT. As a consequence $`\alpha _2`$=1.23 is probably over estimated.
To summarize, the re-brightening by a factor two of the afterglow of GRB 050525a is effective, the transition duration is comprised from less than 3 minutes ($`\delta t/t=0.1`$) up to 17 minutes ($`\delta t/t=0.5`$) centered at $`tt_{\mathrm{trig}}=`$ 32.8 min, and the slope of the temporal decay before and after the re-brightening are fully comparable.
## 3 Early optical GRB afterglows
Only a few afterglows have been observed at optical wavelengths less than one hour after the GRB. They include GRB 990123, GRB 020418, GRB 021004 (z=2.33), GRB 021211 (z=1.01), GRB 041219a, GRB 050319 (z=3.24), GRB 050502a (z=3.79), and GRB 050525a (z=0.606) discussed in this paper. Within this small sample, short-scale variability seems to be the rule rather than the exception. This variability can be observed as single or multiple re-brightenings (GRB 021004, GRB 041219a, GRB 050319, GRB 050525a), as a shallowing (GRB 990123, GRB 021211) or a steepening(GRB 050319) of the light curve, or as a gradual rise of the afterglow (GRB 030418).
Figure 2 compares the early optical afterglows of four GRBs with a measured redshift (GRB 050502a is not included in this comparison since the data on its early afterglow have not yet been published). For a proper comparison, the abscissa of the plot gives the time after the trigger, in the referential of the GRB. A striking feature is the presence of an episode of re-brightening starting 0.01 to 0.02 days after the trigger, in three out of the four GRBs displayed in figure 2.<sup>1</sup><sup>1</sup>1GRB 041219a, whose redshift is not known, exhibits a re-brightening episode starting $``$17 min, or 0.012 day, after the trigger (Blake et al. 2005). At first glance the three GRBs which exhibit rebrightnening episodes do not have special properties. GRB 040319 is a single pulse GRB, while GRB 021004 and GRB 050525a are multi-peak events. The isotropic-equivalent energies and rest frame peak energies of GRBs in figure 2 are 12.6 10<sup>52</sup> erg and 135 keV for GRB 050525a, $`3.1`$ $`10^{52}`$ erg for GRB 050319<sup>2</sup><sup>2</sup>2The peak energy of GRB 050319 is not known., 5.1 10<sup>52</sup> erg and 266 keV for GRB 021004, and 1.4 10<sup>52</sup> erg, and 92 keV for GRB 021211. Since three out of four GRBs with early optical follow-up exhibit re-brightening episodes, it is tempting to conclude that they represent a common feature of GRB afterglows. This remark clearly points out the necessity of very quick optical follow-up to measure the decay slope before 0.01 day (14 min), required to assess the time and amplitude of possible re-brightenings in future GRBs.
## 4 Discussion and conclusion
Within the context of the internal/external shock model of GRBs, re-brightening episodes have been explained by the reverse shock, by a continuing activity of the central engine, by a variable density profile of the external environment in which the fireball expands, by the presence of neutrons in the ejecta, or by the destruction of dust surrounding the source.
Late afterglow emission is explained as the forward shock component of external shock produced by the interaction of the expanding fireball with the external medium. Very early optical emission is thought to be an effect of the reverse shock component of the external shock. The emergence of the forward shock from the reverse shock emission component produces a shallowing of the early-time light curve after an initial steep decay that can mimick a re-brightening (Panaitescu & Kumar panaitescu2004 (2004)). Indeed this model has been recently invoked to explain the phenomenology of GRB 050525a (Shao & Dai Shao05 (2005)).
Variable energy input model also predicts early light curve re-brightenings (Nakar et al. nakar2003 (2003)). The energy variability could be provided by refreshed shocks produced by massive and slow shells ejected late in the GRB, that collide with the inital blast wave when it has decelerated. After the collision of each shell, the flux from the fireball increases but the light curve decay will continue with the same rate as prior to the collision. The net result is a shift upward of the initial light curve at the time of the collision with same decay rate as before collision (Bjornsson et al. bjornsson2004 (2004)). Alternatively, a variable energy input coud be due to initial energy (per solid angle) inhomogeneities in the jet (patchy shell model described by Kumar & Piran 2000b ).
In the variable density scenario, clumpy inter-stellar matter (ISM) or variable wind expelled from the massive progenitor, produces high density regions that, interacting with the expanding fireball, could provide flux enhancements if particular electron cooling conditions are satisifed (e.g. Lazzati et al. lazzati2002 (2002)). In this case, in fact, the flux would be sensitive to density variations only in the ’slow cooling regime’, with the observed frequency below the cooling frequency (Sari et al. Sari98 (1998)). The recovery of the initial slope after the bump could be achieved requiring a decrease in the density below the initial value immediately after the high density region (Lazzati et al. lazzati2002 (2002)).
Afterglow re-brightening has been predicted also by the neutron-fed afterglow model when the already decelerating ion shell sweeps up the trail of decay products left from the ahead decaying neutron shell. The arrival time of the re-brightening depends on the Lorentz factor of the neutron shell and can vary from few seconds to several tens of days after the burst (Beloborodov Beloborodov2003 (2003)).
Finally, dust destruction mechanisms might provide an enhancement of the observed optical flux within a time scale that depends, among several other parameters, on the density distribution of the circumburst dust. A decreasing of reddening is predicted in this case (Perna et al. perna2003 (2003)).
GRB 050525a is, after GRB 050319, the second GRB with a known distance, with detailed observations of the early afterglow at optical (this paper) and X-ray wavelengths (Band et al. band2005 (2005)). This proves that, eigth years after the discovery of GRB afterglows, a new window opens: multi-wavelength observations of the very early afterglow (the first hour). This remarkable advance has been possible thanks to the availability of arcminute localizations quickly distributed to efficient robotic telescopes, and to the excellent performance of the XRT on-board SWIFT. The observations presented in this paper demonstrate the richness of the information contained in the very early afterglow, and the great promises of the multi-wavelength observations which are now within our range.
###### Acknowledgements.
The TAROT telescope has been funded by the Centre National de la Recherche Scientifique (CNRS), Institut National des Sciences de l’Univers (INSU) and the Carlsberg Fundation. It has been built with the support of the Division Technique of INSU. We thank the technical staff joined with the TAROT project: G. Bucholtz, J. Esseric, C. Pollas and Y. Richaud.
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# Explaining Gabriel-Zisman localization to the computer
## 1. Introduction
We formalize the localization of categories, as in the book of Gabriel and Zisman , with the Coq computer proof assistant. The purpose of this preprint is to provide some discussion of this work. On the other hand, the computer files themselves are attached to the companion preprint “Files for Gabriel-Zisman localization”. The text of that preprint consists mainly of the definitions and statements of results from the computer files (in other words it is equal to the files with the proofs removed), plus some instructions for compiling the files.
There are several reasons for choosing this project. A certain amount of basic category theory was done in the files attached to my previous paper on this subject . Thus it is natural to look for some further topics to do in category theory. A long-range goal is to be able to do the theory and practice of closed model categories. A glance at Quillen suggests that the notion of localization of categories à la Gabriel-Zisman is an important component of the statements of some of Quillen’s main results. Also in philosophical terms it is clear that Quillen was influenced by Gabriel and Zisman, so it is reasonable to think that doing a computer formalization of their construction of localization would be a good warm-up exercise.
A little bit of investigation into the bibliographical references for this construction has also turned up another interesting reason to formalize it on the computer. It turns out that the full details of the construction (and specially of the calculus of fractions construction) have never really appeared in print. Or at least, a search for these details has not turned up any reference. Of course it wouldn’t be surprising to find a complete reference somewhere—you might say that this would be the expected normalcy. Nonetheless it seems pretty clear that the vast majority of the very numerous mathematicians who use this theory every day haven’t in fact read a text with the full details written down.
The notion of localization of a category is foundational for some of the most popular tools used by mathematicians today: the homotopy category (of spaces, simplicial sets, or other things), and the derived category (of an abelian category, coming in various flavors). It is surprising that the theory is so hard to find written down in its integrality. This might contribute as part of an explanation for why the theories of homotopy categories and derived categories are so much used and considered as “black boxes”.
One possible point of view would be to say that few have bothered to try to publish the full details of the construction, because in a certain sense that just wouldn’t be worth it: writing something down presupposes that there would be somebody interested in reading it; and writing down the full details of an argument which is in essence straightforward, presupposes that a human reader would desire to, and be capable of, verifying in a meaningful way that the written text really did contain all of the details. Factors such as the total cost of publication also push towards leaving out much of this type of argument.
In trying to write up the present note explaining the computer proof, it became evident that one had to agree with the other published texts on this: the full details of the argument just aren’t sufficiently interesting to justify the rather extensive linguistic effort which would be required to accurately convey them to a human reader, nor interesting enough for the reader to bother reading such an explanation. And this is with a fundamental piece of category theory more than $`40`$ years old. The arrival of the possibility that the “reader” might be a computer changes this calculation. The computer is a perfect listener for an explanation that can be given as a sort of flow of little arguments, sometimes with a necessary global strategy behind them, but always with lots of things to remember, lots of referential notations to refer to various objects, and so forth.
The purpose of the present preprint is to discuss our computer formulation, both of the general localization construction and the special construction when there is a calculus of fractions. We don’t pretend to give all the details in the text—indeed we stop at about the same place as previous authors have. However, the details are necessarily all there in the computer files .
Historically, the notion of localization appeared informally in a somewhat different form in the Tôhoku paper of Grothendieck where he formalizes (to some extent) language which he attributes to Serre of working in a category “modulo” a subcategory.
After Gabriel-Zisman, the question of localization of categories has been treated in a number of references. Several people were helpful in pointing out some of these in response to requests posted to the topology and category-theory mailing lists. The references include books by L. and N. Popescu , H. Schubert , and F. Borceux . Curiously enough the localization construction doesn’t seem to appear in , although the underlying free and quotient category techniques are there. A classical reference which chronologically goes alongside Gabriel and Zisman is Verdier’s thesis but which was only recently published . Verdier considers the case of localization of an additive category, inverting a multiplicative system which satisfies a two-sided calculus of fractions. The construction is similar to the left-fraction construction.
The above list of references is undoubtedly partial. It doesn’t include very many of what are certainly numerous research papers since the time of which may treat aspects of these issues in some detail ( is an example).
Nonetheless, it is interesting to note the prevalence of formulations leaving “to the reader” parts of the proofs of details of the localization constructions. For example (the following all refer to the left or right fractions construction):
, p. 117: “La preuve est facile et laissée au lecteur qui pourra démontrer de même la proposition ci-après …”.
, p. 260: “Using (i), (ii), (iii) and (v), there is no difficulty in verifying that the composite …is well defined, …”.
, p. 155: “It is not difficult to see that with the equivalence relation (2) introduced above one also has a well-defined composition law …”.
Another interesting reference is Pronk’s paper on localization of $`2`$-categories , pointed out to me by I. Moerdijk. This paper constructs the localization of a $`2`$-category by a subset of $`1`$-morphisms satisfying a generalization of the right fraction condition. In this case, there is no need to divide by an equivalence relation on the set of $`1`$-arrows, because the appropriate arrows (i.e. pairs or what we call “fraction symbols”) are identified by the presence of $`2`$-cells making them equivalent. On the other hand, knowing which $`2`$-cells are there would tell us which fraction-symbols need to be identified in the $`1`$-localization. Thus it seems likely that from the high level of detail present in , one could extract the complete set of necessary arguments for the localization of a $`1`$-category. Nonetheless, the full set of details for the coherence relations on the level of $`2`$-cells is still too much, so the paper ends with:
, p. 302: “It is left to the reader to verify that the above defined isomorphisms $`a`$, $`l`$ and $`r`$ are natural in their arguments and satisfy the identity coherence axioms.”
We take the opportunity at this point in the introduction to mention the colimit point of view about localization, even though it isn’t treated in our proof verification files. The contents of this discussion are touched upon by Gabriel-Zisman in 1.5.4 of Chapter I, see also 6.2 of Chapter II, and since then has become even more well-known.
The localization $`𝒞[\mathrm{\Sigma }^1]`$ can be viewed as a pushout or colimit in the category of categories (one has to fix a universe $`𝒰`$ and consider colimits in the category of $`𝒰`$-categories). To be precise, it is a pushout fitting into a cocartesian diagram
$$\begin{array}{ccc}\mathrm{\Sigma }\times I& & 𝒞\\ & & \\ \mathrm{\Sigma }\times \overline{I}& & 𝒞[\mathrm{\Sigma }^1]\end{array}$$
where here $`\mathrm{\Sigma }`$ denotes the discrete set of maps to be inverted (considered as a discrete category), $`I`$ denotes the category with objects $`0`$ and $`1`$ and one non-identity arrow $`01`$, and $`I\overline{I}`$ denotes the completion to a category with two objects $`0`$ and $`1`$ joined by a single isomorphism. In a certain sense the fact that this diagram is cocartesian just restates the universal property of the localization.
Ross Street in mentionned the above pushout point of view as well as another closely related terminology, saying that the localization is the “coinverter” of the 2-cell $`\sigma `$ which is the natural transformation from
$$\mathrm{dom}:\mathrm{\Sigma }C\text{to}\mathrm{cod}:\mathrm{\Sigma }C.$$
The fact that the localization is a pushout implies that this operation is compatible with colimits of categories. For example suppose
$$\begin{array}{ccc}A& & B\\ & & \\ C& & P\end{array}$$
is a cocartesian diagram of categories, and $`\mathrm{\Sigma }_A`$, $`\mathrm{\Sigma }_B`$, and $`\mathrm{\Sigma }_C`$ are subsets of morphisms in $`A`$, $`B`$ and $`C`$ respectively such that $`\mathrm{\Sigma }_A`$ maps into $`\mathrm{\Sigma }_B`$ and $`\mathrm{\Sigma }_C`$. Let $`\mathrm{\Sigma }_P`$ be the union of the images of $`\mathrm{\Sigma }_B`$ and $`\mathrm{\Sigma }_C`$. Then the diagram
$$\begin{array}{ccc}A[\mathrm{\Sigma }_A^1]& & B[\mathrm{\Sigma }_B^1]\\ & & \\ C[\mathrm{\Sigma }_C^1]& & P[\mathrm{\Sigma }_P^1]\end{array}$$
is cocartesian. Let’s stress that none of this is proven in the proof files.
An interesting case is when we take $`\mathrm{\Sigma }=Mor(C)`$ to be the full set of morphisms of a category. In this case the localization is the groupoid-completion, the universal groupoid with a functor from $`C`$, denoted
$$C^{\mathrm{gr}}:=C[Mor(C)^1].$$
The compatibility with colimits stated above gives the following statement for groupoid completions: if
$$\begin{array}{ccc}A& & B\\ & & \\ C& & P\end{array}$$
is a cocartesian diagram of categories then the diagram of groupoids
$$\begin{array}{ccc}A^{\mathrm{gr}}& & B^{\mathrm{gr}}\\ & & \\ C^{\mathrm{gr}}& & P^{\mathrm{gr}}\end{array}$$
is cocartesian (either in the category of $`𝒰`$-categories, or in the category of $`𝒰`$-groupoids). The groupoid completion $`C^{\mathrm{gr}}`$ is equivalent to the Poincaré fundamental groupoid of the realization of the nerve of $`C`$ (denoted by $`|C|`$ for short). In view of which, the above statement can be viewed as a sort of Van Kampen theorem for fundamental groupoids in the style of R. Brown . To make this view totally precise we would have to look at when the corresponding diagram of spaces $`|A|`$, $`|B|`$, $`|C|`$, $`|P|`$ is a pushout of spaces, which is sometimes but not always the case.
Ross Street also pointed out to me in , that the problem of localizing a noncommutative ring is closely related (and somewhat similar to the Verdier case in that it brings in additive structure). He sent me some notes treating in great detail the ring localization; it should be relatively easy to transform those into a computer proof for this case. On the subject of notes, Clark Barwick mentions that he had written up some notes about Lemma 1.2 of Gabriel-Zisman; there are probably (one would hope) a number of mathematicians who have done so.
Here is the plan of the paper. We will discuss first the general construction of the localization in a mathematical fashion; this is followed by a section discussing the issues which arise in the computer formulation. Then we come back to a mathematical discussion of the calculus-of-fractions construction, including discussion of the subtleties which arise when we go towards the full details of the argument. The section after that discusses the new issues which arise in the computer formulation, notably how we deal with the commutative diagrams which one is tempted to use for the proof. In the last section, we mention very briefly the contents of the remaining files in the present development.
## 2. The general localization construction
We recall in usual mathematical terms how the general construction of the localization of a category works, taken directly from the first two pages after the introduction in Gabriel-Zisman . Fix a category $`𝒞`$ and a subset $`\mathrm{\Sigma }Mor(𝒞)`$ of morphisms. We make no assumption about $`\mathrm{\Sigma }`$. We construct the localization, denoted $`𝒞[\mathrm{\Sigma }^1]`$, as follows. Start by taking the disjoint union $`Mor(𝒞)\mathrm{\Sigma }`$. The arrows in the first factor are thought of as going in the forward direction, and the arrows in the second factor as going in the backward direction. This allows us to create a directed graph whose vertices are the objects of $`𝒞`$ and whose set of edges is $`Mor(𝒞)\mathrm{\Sigma }`$. Let $``$ denote the free category over this graph. Recall that this means that the objects of $``$ are the vertices of the graph (thus, the same as the objects of $`𝒞`$), and the morphisms of $``$ are directed paths in the graph.
Next introduce some relations on $``$, and let $`𝒞[\mathrm{\Sigma }^1]`$, which we denote as $``$ for short, be the quotient of $``$ by these relations. The relations are trivial on the set of objects, that is to say no different objects are put in relation and related morphisms share the same source and target. The relations are introduced with the purpose of insuring the following properties for the quotient $``$:
(1) the natural map on arrows $`Mor(𝒞)Mor()`$ (which is not itself a functor) projects to a functor $`𝒞`$; and
(2) if $`u\mathrm{\Sigma }`$ then the image in $`𝒞`$ of the backward edge corresponding to $`u`$ (that is, the element of the second factor of the disjoint union) is inverse in $``$ to the image of $`u`$ by the functor in (1).
The relations are chosen heuristically in a minimal way to accomplish this. In Gabriel-Zisman this process is written in a compact way: the very process of taking the quotient category implies that we want the set-theoretical quotient of $`Mor()`$ by the full set of relations, to be the set of morphisms of a category. This in itself contains some properties of compatibility between the relations and the composition and identity operations. Call a relation which satisfies these properties, a categorical relation. We can start by specifying an arbitrary list of relations and then take the closure under this condition, that is the smallest categorical relation containing our list of relations. Once we have decided to do this, we can list the germinal relations as follows: condition (1) requires that we specify two types of relations, one for compatibility of the functor with the composition, and one for the compatibility of the functor with the identity; and condition (2) requires again two types of relations, one each for the left and right inverse properties. Thus we need to impose $`4`$ families of relations to start with. These are listed on page 6 of Gabriel-Zisman (as well as in our proof file, see below). Given this list of relations, the quotient $``$ is constructed by first completing to the smallest categorical relation $``$ containing the list, then taking $`Mor():=Mor()/`$. It is straightforward to show that this defines a category $``$, with a functor $``$ satisfying properties (1) and (2) above, since we constructed the relation that way on purpose. We obtain a functor $`P_\mathrm{\Sigma }:𝒞`$ sending elements of $`\mathrm{\Sigma }`$ to invertible morphisms in $``$.
So much for the construction. The next step is to state and prove an appropriate universal property. On the first page of the construction they give the main properties which are (I am quoting):
“(i) $`P_\mathrm{\Sigma }`$ makes the morphisms of $`\mathrm{\Sigma }`$ invertible,
(ii) If a functor $`F:𝒞𝒳`$ makes the morphisms of $`\mathrm{\Sigma }`$ invertible, there exists one and only one functor $`G:𝒞[\mathrm{\Sigma }^1]𝒳`$ such that $`F=GP_\mathrm{\Sigma }`$.”
After explaining the construction in $`14`$ lines at the bottom of page 6 and the top of page $`7`$ (which I have paraphrased above), Gabriel-Zisman jumped right up to a highly abstract formulation of the universal property which we recopy here together with a few subsequent phrases:
1.2. Lemma: For each category $`X`$, the functor $`\underset{¯}{Hom}(P_\mathrm{\Sigma },X):\underset{¯}{Hom}(𝒞[\mathrm{\Sigma }^1],𝒳)\underset{¯}{Hom}(𝒞,𝒳)`$ is an isomorphism from $`\underset{¯}{Hom}(𝒞[\mathrm{\Sigma }^1],𝒳)`$ onto the full subcategory of $`\underset{¯}{Hom}(𝒞,𝒳)`$ whose objects are the functors $`F:𝒞𝒳`$ which make all the morphisms of $`\mathrm{\Sigma }`$ invertible.
“The proof is left to the reader. This lemma states more precisely conditions (i) and (ii). From now on …”
Afterwards they pass immediately to the discussion of motivating examples like when the multiplicative system comes from a pair of adjoint functors.
This text (which totals less than a full page) is quite interesting from the point of view of the problem of formalizing mathematics on the computer. In a very short space the authors have indicated, without error and indeed giving all of the necessary information, a relatively complex mathematical construction, together with a very abstract statement of the universal property satisfied by this construction. Starting with the information given here, it is a straightforward (and mathematically uninteresting) exercise to fill in all of the required details.
Creating a proof document to be read by a computer proof assistant raises a certain number of mild difficulties. As a test, I have tried to attain the exact statement of Lemma 1.2 given above. Before getting to a discussion of some details of this process in the next section, it is interesting to note here that the time it took to do this was about a month, and the resulting total size of the $`3`$ proof files involved (`freecat.v`, `qcat.v` and `gzdef.v`) is about 10,000 lines.
## 3. The computer formulation
In order to formalize the general localization construction, we use the Coq proof assistant . We place ourselves in an axiomatic environment which implements classical Zermelo-Fraenkel set theory, and also relates it to the type theory of Coq. Concretely, the proof files attached to include identical copies of the proof files of my earlier set theory and category theory developments . The set-theoretical part of this starts with a file axioms.v containing all of the axioms we assume (that is, no subsequent files use the Axiom or Parameter commands). These axioms are intended to implement ZFC within Coq. However, we furnish no formal proof of the fact that they do indeed do that, in other words that together with the Coq “calculus of inductive constructions” type system, they furnish a mathematical system which is consistent within the context of the usual ZFC axioms. It would be good to have such a proof, but that seems to be complicated (due to the complicated structure of Coq) and possibly nontrivial due to certain aspects of Coq’s type system such as cumulativity between sorts Prop and Type. One would have to prove Conjecture 2 of Miquel and Werner . This is left to the reader!
The next question which is left open is to convince oneself that the definitions and statements of lemmas contained both in the present files as well as in the category theory files, accurately represent what the mathematician means when he speaks of categories, functors and so forth. This again may contain some nontrivial aspects and is left to the reader. The accompanying preprint “Files for Gabriel-Zisman Localization” is intended to help with this task: the textual part of this preprint consists of the Coq files, with all proofs taken out. There one can look directly at the definitions and statements of lemmas, which are the only parts which need to be understood in order to verify the meaning of what is being said. (However, this is done only for the files concerning localization; in principle it might be a good idea to have the same thing for the set-theory and category-theory developments but that would be lengthy.)
While speaking about these foundational questions for the computer formulation, it is important to note that one could undoubtedly use any of a number of other environments for treating this question. For example, it should be possible to proceed based on Saibi’s category theory contribution where sets are replaced by “setoids” (types plus equivalence relations). This is particularly so in that one major element of the localization construction is the notion of quotient category. It is likely that a setoid approach would simplify certain aspects, at the price of introducing other complications elsewhere. We don’t venture to predict how economical that would be on the whole. One should of course also envision doing this type of formalization within other proof assistants (the list of which is getting very long and we don’t attempt to reproduce it here, see ).
For the reader who, at this point, still feels that a computer formalization can add something to the question of verifying the mathematics underlying the localization construction, we now consider some details.
### 3.1. Category theory
We only treat small categories, i.e. ones whose objects and morphisms form sets. In this context any distinction between small and big categories would be made by refering to a Grothendieck universe (which would itself be a set). However, in the current development we don’t treat the question of when the localization of a category which is big but has small $`hom`$ sets with respect to a given universe, is again big but with small $`hom`$ sets with respect to that universe. Thus, we are always working with sets and no foundational acrobatics come into play.
Most major notational questions about categories have already been dealt with in the category theory files. A category is a $`5`$-uple consisting of the set of objects, the set of morphisms, the graph of the partially defined composition operation, the graph of the identity operation, and a fifth set which is destined to contain any eventual extra structure one might want to include. The positions in the $`5`$-uple are indicated by character strings (i.e. the $`5`$-uple is a function whose domain is a set of $`5`$ specific character strings corresponding to the $`5`$ places). This schema isn’t the most economical: the necessary data (excepting the last structure variable) is contained in the graph of the composition morphism. The goal is rather to achieve some rudimentary standardization of the procedure for considering mathematical objects.
One feature of the category-theory encoding which is worthwhile to recall here is that the set of morphisms is supposed to contain only objects `u` which themselves are triples containing a source, a target, and a third indicative element. This property is written `Arrow.like u`. Here as before, these triples are realized as functions whose domain is a set of $`3`$ character strings. This allows us to consider `source u` and `target u` for an arrow `u`, independantly of the category for which `u` is a morphism. Here we have an economy of notation which has been extremely useful throughout the category theory development. In our discussion below we will encounter several places where a certain modification of the “obvious” approach is made necessary by the `Arrow.like` hypothesis. These modifications are easy to do once we are aware of the phenomenon—which is why I am devoting some space below to these explanations.
Similar notational considerations hold for functors and natural transformations. We refer to for further discussion of these issues.
### 3.2. The free category
The first step in our current files is `freecat.v` where we construct the free category on a graph. Since a morphism in the free category is a path in the graph, we need to implement the notion of path. This touches on what G. Gonthier explained was an important piece of their work formalizing the 4-color theorem . However, the approaches are not the same since we are much less concerned with efficient computation on these objects and more concerned with their theoretical manipulation. The notion of path also appears in T. Hales’ recent formalization of the Jordan curve theorem .
To implement a notion of “path”, we obviously need a theory of “uples”, which are implemented as functions whose domain is an interval of natural numbers of the form $`[0,\mathrm{},n1]`$ where $`n`$ is the length of the uple. We define the function `Uple.create` to create an uple of length $`l`$ from a function `f:nat -> E`, and a function `component` to get back the $`i`$th element of an uple. We need the `length` function as well as `concatenate`.
An important lemma is `uple_extensionality` which says that two uples of the same length with the same elements are the same; this allows us to prove associativity of concatenation. J.S. Moore said for his ACL2 system that the first thing you would want to prove was associativity of concatenation. In that type of system, uples or lists are inductive objects and associativity is a statement proved by recurrence on the length. In Gonthier’s paper , the notion of path is defined structurally so that associativity is automatic by term reduction and need not be mentionned as a lemma. In our case, the technical tool used to simplify the proof of associativity, and most other manipulations of our uples which are functions of natural numbers, is the `omega` tactic. The usefulness of this tactic, developped by Crégut based on an algorithm of Pugh , was pointed out to me by Marco Maggesi. It dispatches easily any arithmetic statement involving the standard operations and inequalities on natural numbers. In our proof files this powerful tactic is abbreviated as `om` which could be interpreted alternatively as a reference to Buddhism or the Marseille soccer team. Finishing the `Uple` module is the operation `utack` which corresponds to concatenation with an uple of length one. This specific case enters often later so we treat it specifically.
The notion of graph is relatively easy to encode. A graph is a pair consisting of a set of vertices, and a set of edges. The edges of a graph are also supposed to be arrows. This situation is simple enough to provide a good example of our general notational procedure which we can recopy here:
```
Definition Vertices := R (v_(r_(t_ DOT ))).
Definition Edges := R (e_(d_(g_ DOT))).
Definition vertices a := V Vertices a.
Definition edges a := V Edges a.
Definition create v e :=
denote Vertices v
(denote Edges e stop).
Definition Graph.like a := a = create (vertices a) (edges a).
Definition Graph.axioms a := Graph.like a &
(forall u, inc u (edges a) -> Arrow.like u) &
(forall u, inc u (edges a) -> inc (source u) (vertices a)) &
(forall u, inc u (edges a) -> inc (target u) (vertices a)).
```
We don’t do any theory of graphs beyond just the definition.
Next we look at the paths which will make up the morphisms of the free category on a graph. These are arrows whose third term are uples; and furthermore the uples will eventually (in the definitions `arrow_chain` and `mor_freecat`) be supposed to be sequences of composable arrows in the graph, starting from the source of the arrow and ending at the target. This situation requires a certain amount of specific treatment, for example we define a version `segment` of the previous `component` function (and `seg_length` instead of `length`). In general terms, this type of definition contracting two or some other small number of functions which often occur together, necessitating the transposition of all of the lemmas concerning the pieces, occurs all over the place and seems to be a general phenomenon. The composition operation for the free category is defined by using concatenation of the underlying uples. We also define the identity (whose uple has length $`0`$) and prove all the various things needed to obtain the category axioms. We then would like to consider functors from the free category into another category. For this, we need to define the operation of composing together a composable sequence of arrows in a category (the definition `mor_chain` is very much like `arrow_chain`). In this way we can state a universal property of the free category on a graph (see the results concerning the construction `free_functor`). To close out this discussion we also consider (in the results concerning `free_nt`) natural transformations between functors whose sources are the free category. This is significantly easier because a natural transformation is a function on objects, and the objects of the free category are just the vertices of the graph.
### 3.3. Quotient categories
After the free category, the other main element of the construction is the notion of quotient category (`qcat.v`). We are in a situation which is significantly easier than the general case: our relation has no effect on the objects. In other words, we have a relation on the set of morphisms of a category, such that two morphisms which are related already have the same source and target. It is convenient to distinguish two separate notions, denoted `(cat_rel a r)` and `(cat_equiv_rel a r)`. The first means that `r` is an arbitrary relation on the morphisms of a category `a`, respecting source and target. The second means that `r` is an equivalence relation and compatible with the composition of `a`. One important example of a `cat_equiv_rel` is `(coarse a)` which puts in relation any two morphisms with the same source and target. The existence of this maximal relation allows us by intersection to define the smallest `cat_equiv_rel` containing a given `cat_rel` `r`. We call this construction `(cer a r)` (here `cer` stands for the “`c`ategorical `e`quivalence `r`elation” on the category `a` generated by `r`).
In order to construct the quotient category of `a` by `r`, we need a manipulation called `arrow_class`. The reason for this is that in our notion of category, the morphisms are supposed to be `Arrow.like`, i.e. triples having a source, a target and an arrow. Thus we can’t just say that the set of morphisms of the quotient category is the usual quotient (i.e. set of equivalence classes) of the set of morphisms by the relation. Thus we define `(arrow_class r u)` to be the arrow with the same source and target as `u`, but whose third element is the equivalence class of `u` for the relation `r`. Now the set of morphisms of the quotient category will be the image of this construction as `u` runs through the morphisms of `a`. The definition `(is_quotient_arrow a r u)` formalizes the statement that `u` is in the image. We also need a construction `(arrow_rep v)` going in the other direction (see Lemmas `related_arrow_rep_arrow_class` and `arrow_class_arrow_rep` saying that the two constructions are inverse in the appropriate sense). We then define `quot_id` and `quot_comp`, the operations which will become the identity and composition for the quotient category. As usual, before trying to construct the category it is a good idea to prove all of the necessary properties for these constructions. Then when we construct `(quotient_cat a r)` we prove destruct-create lemmas `comp_quotient_cat` and `id_quotient_cat` saying that the identity and composition are `quot_id` and `quot_comp`. The destruct-create lemmas `ob_quotient_cat` and `mor_quotient_cat` are proven after `quotient_cat_axioms` because the properties `ob` and `mor` include the category axioms for their first variables.
The module `Quotient_Functor` does similar things for defining a functor `qfunctor` to the quotient category, and a functor `qdotted` from the quotient category. The latter terminology is intended to suggest that `qdotted` is the dotted arrow which is filled in in the universal property of the quotient category. Thus if `f` is a functor and `r` a categorical equivalence relation on the category `source f` we get a functor `qdotted r f` such that
```
source (qdotted r f) = quotient_cat (source f) r
target (qdotted r f) = target f
fcompose (qdotted r f) (qprojection (source f) r) = f.
```
The unicity statement for the universal property says that if `f` is a functor with
```
source f = quotient_cat a r
```
then
```
f = qdotted r (fcompose f (qprojection a r)).
```
There are no particular difficulties encountered in these arguments beyond the kind we have already discussed above.
Also contained in the file `qcat.v` is a module `Ob_Iso_Functor` dedicated to studying the following situation. We have a functor `f` and a category `a`. We study the pullback morphism induced by `f`, denoted `(pull_morphism a f)`, from `(functor_cat (target f) a)` to `(functor_cat (source f) a)`. Recall that these constructions come from the file on functor categories `functor_cat.v` in the category-theory development. The purpose of this module is to contribute to the proof of Gabriel-Zisman’s Lemma 1.2. In particular, we will want to apply this to the case where `f` is the functor from a category to its localization. Thus we assume that `f` is an isomorphism on objects. We develop a criterion for when `(pull_morphism a f)` is fully faithful and injective on objects (see the definition `iso_to_full_subcategory`), or equivalently that it induces an isomorphism from `(functor_cat (target f) a)` to a full subcategory of `(functor_cat (source f) a)`. The equivalence between these notions is shown in Lemma `iso_to_full_subcategory_interp`.
Intervening in the statement of the criterion is the construction `add_inverses a s`. This is the subset of morphisms of `a` which are either already in `s`, or else are inverses in `a` to morphisms in `s`. Our criterion, stated at the end of this module in Lemma `iso_to_full_subcategory_pull_morphism_criterion`, says that if a functor `f` is an isomorphism on objects, and if `(add_inverses (target f) (mor_image f))` generates the category `(target f)`, then for any category `a` the pullback functor `pull_morphism a f` is an isomorphism onto a full subcategory. We will use this criterion, applied to the functor from a category to its localization, to obtain half of the statement of Lemma 1.2.
Finishing out the file `qcat.v` is a module `Associating_Quotient` which substantially recopies much of the definition of quotient category. The only difference is that we don’t start with a category but only with a structure like a category but which doesn’t necessarily satisfy the associativity or left and right identity axioms. The idea is that the equivalence relation will enforce these axioms. This construction is not needed for the general construction of localization, but it will be needed later for the construction of the category of fractions. It didn’t seem necessary to go back and redo the whole quotient construction with this generality in mind: it is easier to recopy the relevant parts and change them. This might result in a file which is longer than necessary, but one should keep in mind that the variable we are trying to economize is the energy necessary to produce (or understand) the collection of files, not their total length.
### 3.4. Construction of the localization
Recall that the construction of the localization starts by looking at the graph whose edge set is the disjoint union of the morphisms of `a` with the elements of `s`. Since these two sets are anything but disjoint, we need some additional notation to implement the disjoint union. To this end, the first thing one notices at the start of the file `gzdef.v` is the introduction of two sets `Forward` and `Backward`. These are character strings (which are elements of `E` hence sets, see `notation.v` ). An edge of the graph is either a “forward edge” corresponding to a morphism in `a`, or else a “backward edge” corresponding to an element of the localizing system `s`. The obvious thing is to try putting
```
forward_arrow u := pair Forward u
backward_arrow u := pair Backward u.
```
However, this doesn’t work. The reason is that the elements of the set of edges of the graph are supposed to be `Arrow.like`. To remedy this problem, we set
```
forward_arrow u := Arrow.create (source u) (target u) (pair Forward u)
backward_arrow u := Arrow.create (target u) (source u) (pair Backward u).
```
Notice that the source and target are interchanged in the function `backward_arrow`. The function `original_arrow` yields back the arrow we started with. Now `loc_edges a s` is the set of such edges, i.e. the union of the images of `forward_arrow` and `backward_arrow` respectively on the morphisms of `a` and on `s`. The union is disjoint because `Forward` and `Backward` are distinct. Define the graph `gz_graph a s` whose vertices are the objects of `a` and whose edges are `loc_edges a s`.
From here, the construction basically follows the ordinary one, and doesn’t really take up too much space. The free category on `gz_graph a s` is called `gz_freecat a s`. The definition `gz_rel a s` is where the defining relations for the construction of the localization are listed. We recopy here a lemma which rewrites that definition in a slightly more readable fashion.
```
Lemma related_gz_rel : forall a s e f, localizing_system a s ->
related (gz_rel a s) e f =
((exists x, (ob a x &
e = (forward_edge (id a x)) &
f = (freecat_id x))) \/
(exists q, (inc q s &
e = (freecat_comp (forward_edge q) (backward_edge q)) &
f = (freecat_id (target q)))) \/
(exists q, (inc q s &
e = (freecat_comp (backward_edge q) (forward_edge q)) &
f = (freecat_id (source q)))) \/
(exists u, exists v, (mor a u & mor a v & source u = target v &
e = (freecat_comp (forward_edge u) (forward_edge v)) &
f = (forward_edge (comp a u v))))).
```
The size of this text is comparable to the size of the paragraph of where the relations are listed. Then `gz_cer a s` is the associated categorical equivalence relation, and `gz_loc a s` is the quotient category. The functor from `a` to `gz_loc a s` is called `gz_proj a s`.
The module `GZ_Thm` is where we prove Gabriel-Zisman’s Lemma 1.2. From the file `qcat.v`, the module `Ob_Iso_Functor` furnishes the results necessary to prove the part of Lemma 1.2 which says that pullback is an isomorphism onto a full subcategory. As pointed out above, one of the delicate points is that the statement of Lemma 1.2 involves the pullback morphism `pull_morphism` between functor categories. Functor categories were treated in the category theory development , and the place where we made use of that theory was in the module `Ob_Iso_Functor` so we don’t actually encounter it too much anymore here. This conclusion is stated in the present file as `iso_to_subcategory_pull_gz_proj`, a corollary of the fact that `gz_loc` is generated by adding available inverses to the morphism image of the functor `gz_proj`.
The main part of the work done in the present module is to prove the versal part of the universal property. Furthermore, rather than just giving a proof we would like to have some useful notation. We start by introducing this notation for the free category: the operation `fr_dotted` corresponds to filling in the dotted line in a diagram expressing the versality of the universal property. Similarly `qdotted` did the same thing in the file `qcat.v`, and putting them together we get a construction called `gz_dotted` which expresses versality in the following way. Given `a`, `s` and a functor `f` with `source f = a`, we say `loc_compatible a s f` if `f` sends elements of `s` to invertible morphisms in `target f`. Then `gz_dotted a s f` is a functor with
```
source (gz_dotted a s f) = gz_loc a s
target (gz_dotted a s f) = target f
fcompose (gz_dotted a s f) (gz_proj a s) = f
```
The last property here, which is Lemma `fcompose_gz_dotted_gz_proj`, corresponds to the versality property (ii) of , page 6. The uniqueness property (i) on page was our Lemma `gz_proj_epimorphic`. These are actually the properties which are the most useful in practice.
Our version of Lemma 1.2 is given by two statements,`iso_to_subcategory_pull_gz_proj` as noted above, and for identification of the full subcategory image of `pull_gz_proj`, the lemma `ob_image_pull_gz_proj`.
It might eventually be useful to have a more concrete description of natural transformations between functors starting from the localization, but apart from the fact that it is implicitly contained in the statement of Lemma 1.2, we don’t treat this further here.
## 4. Calculus of left (or right) fractions
When I gave a talk in Nice about the computer formulation of the general localization construction, Charles Walter suggested that it would be interesting to compare the formalization of the general construction of localization, with what would have to be done to construct the localization in the presence of the habitual calculus of fractions conditions. With this motivation I set out a while later to formalize the fractions construction from Chapter 2 of Gabriel-Zisman.
Contrarily to the general construction, it turned out (in my own opinion at least) that filling in the details of the left-fractions construction involved some nontrivial (if easy) mathematical thought, and drawing lots of diagrams. We don’t draw diagrams in the computer formulation (that might someday be possible but it is beyond the reach of most computer proof assistants for the moment). As a replacement, we set up definitions of situations involving several arrows of a category, which correspond to the diagrams we would want to draw. This will be discussed in the next section.
In the present section we go into some detail about the mathematics of the problem, which stems from the fact that Gabriel-Zisman state their fraction construction under a somewhat weak collection of hypotheses about the localizing system.
The dual notions of left and right calculus of fractions are intended to be analogues of the notion of multiplicative system for a commutative ring, which as was well-known leads to a description of the localization as a set of “fractions”. In the case of categories, one would like to represent elements of the localizations as “fractions” or diagrams
$$x\stackrel{v}{}y\stackrel{t}{}z,$$
where by convention the arrows going backward are supposed to be in $`\mathrm{\Sigma }`$. This diagram is viewed as representing the morphism $`t^1v`$ of $`𝒞[\mathrm{\Sigma }^1]`$ so it is called a left fraction symbol. We would like to have a nice set of conditions guaranteeing first of all that every morphism of the localization can be written as a fraction; and second guaranteeing that the equivalence relation on formal symbols $`(t,v)`$ whose quotient the set of morphisms $`t^1v`$ is easy to understand. This collection of conditions is the calculus of left fractions. There will be a dual notion of calculus of right fractions obtained by conjugating everything with the ‘opposite’ construction. Aside from the problem of implementing this conjugation in the computer formulation, we will focus on the left-fraction case.
The conditions for a calculus of left fractions are given on page 12, (2.2 a,b,c,d). For convenience we reproduce them here:
(a) $`\mathrm{\Sigma }`$ contains the identity morphisms of all objects of $`𝒞`$;
(b) $`\mathrm{\Sigma }`$ is closed under composition;
(c) If
$$X^{}\stackrel{s}{}X\stackrel{u}{}Y$$
is a diagram with $`s\mathrm{\Sigma }`$ then there exists a commutative square
$$\begin{array}{ccc}X& & Y\\ & & \\ X^{}& & Y^{}\end{array}$$
such that the right downward map in the square is in $`\mathrm{\Sigma }`$; and
(d) If $`f`$ and $`g`$ are two maps from $`X`$ to $`Y`$ such that there exists $`s\mathrm{\Sigma }`$ with $`fs=gs`$, then there is a morphism $`t:YY^{}`$ in $`\mathrm{\Sigma }`$ such that $`tf=tg`$.
Condition (c) says that every right fraction symbol can be completed to a left fraction symbol (in the commutative square, $`X^{}Y^{}Y`$ is a left fraction), that is dividing by an element of $`\mathrm{\Sigma }`$ on the right can be changed to division on the left. Condition (d) says that equalization on the right can be changed to equalization on the left.
Most notable about this definition, specially in light of common practice in more recent times, is what is left out. It is natural to require the following condition, which we call three for two:
(e) if $`X\stackrel{g}{}Y\stackrel{f}{}Z`$ is a composable pair of morphisms, then if any two of $`f`$, $`g`$ and $`fg`$ are in $`\mathrm{\Sigma }`$, the third one is too.
In general we can define the saturation of a set of morphisms to be the set $`\mathrm{\Sigma }^{\mathrm{sat}}`$ of all morphisms in $`𝒞`$ which become invertible in $`𝒞[\mathrm{\Sigma }^1]`$. It is clear from the universal property that the functor
$$𝒞[\mathrm{\Sigma }^1]𝒞[(\mathrm{\Sigma }^{\mathrm{sat}})^1]$$
is an isomorphism, that is a set $`\mathrm{\Sigma }`$ and its saturation share the same localization. It is also clear that for any set of morphisms, the saturation satisfies conditions (a), (b) and (e). Thus from a certain perspective there would be no loss of generality in requiring that our set of morphisms satisfy condition (e). For example, Quillen will later incorporate this condition as an important part of his notion of “closed model category”.
Nonetheless, Gabriel-Zisman don’t make this requirement (and indeed they don’t even speak of the three-for-two condition (e) near here in the text). When you start to look closely at the details it becomes clear that stating and proving the construction of the left-fraction localization in the absence of the three-for-two condition is a bit of a challenge, one which they happily ask the reader to meet almost without saying anything about it, just subtlely giving the correct definition of the equivalence relation so as to make it work.
Throughout the discussion of the left-fraction condition—as was the case for the general construction too—Gabriel-Zisman make reference to the construction of the sets of morphisms as being a direct limit construction. We ignore this aspect here: it isn’t treated in the formal proof development and we don’t discuss it in the informal presentation either. In fact it goes beyond the concrete character of the construction and it isn’t clear whether it represents a useable piece of information (although that doesn’t mean that it isn’t conceptually important).
We now get to the description of the equivalence relations. We define a preliminary set of formal symbols $`(t,f)`$ consisting of two arrows having the same target, the first of which is in $`\mathrm{\Sigma }`$. A left fraction symbol $`(t,f)`$ is drawn as a diagram
$$x\stackrel{f}{}y\stackrel{t}{}z.$$
We would like to define the set of morphisms of the left-fraction category to be the quotient of this preliminary set by an equivalence relation (, the top of page 13). Before stating the relation, notice that the “source” of the formal symbol $`(t,f)`$ is the source of $`f`$, whereas the “target” of $`(t,f)`$ is defined to be the source of $`t`$. We call the common target of $`t`$ and $`f`$ the vertex of the symbol. The equivalence relation will preserve source and target. Two symbols $`(s,f)`$ and $`(t,g)`$ are said to be equivalent if there are maps $`a`$ and $`b`$ such that the source of $`a`$ is the vertex of $`(s,f)`$ and the source of $`b`$ is the vertex of $`(t,g)`$, and $`af=bg`$, $`as=bt`$, and furthermore $`as=bt`$ is in $`\mathrm{\Sigma }`$. Note that these conditions automatically say that the targets of $`a`$ and $`b`$ are the same. See the second diagram on page 13 of Gabriel-Zisman.
We can think of these conditions as giving a symbol $`(as,af)=(bt,bg)`$ which is “beyond” both $`(s,f)`$ and $`(t,g)`$, and indeed this notion is what we use in the proof development. We say that a symbol $`(r,u)`$ is beyond $`(s,f)`$ if there exists a morphism $`a`$ whose source is the vertex of $`(s,f)`$ and target the vertex of $`(r,u)`$, such that $`r=as`$ and $`u=af`$. In this case we say that the morphism $`a`$ is an intermediary from $`(s,f)`$ to $`(r,u)`$.
A natural impulse would be to ask that the morphism $`a`$ (or the morphisms $`a`$ and $`b`$ in the definition of the equivalence relation) be in $`\mathrm{\Sigma }`$. This would be automatic from the conditions that $`s`$ and $`as`$ are in $`\mathrm{\Sigma }`$, if we had the three-for-two condition (e). However, if we try to do the construction in the absence of (e), we shouldn’t ask that the intermediary morphism $`a`$ be in $`\mathrm{\Sigma }`$ because then the construction wouldn’t work. A counterexample is discussed in the file `lfcx.v`.
If we have condition (e), then the proof that this defines an equivalence relation is relatively straightforward. Without it, things are somewhat more tricky. The details necessary to overcome this problem must be considered as subsumed in the phrase “It follows from (a), (b), (c), (d) that this defines an equivalence relation …” in the middle of page 13 . We will now explain how to see that.
The main difficulty lies in proving that the equivalence relation is transitive. This may be rewritten in terms of the notion of “beyond” as trying to show<sup>∗\*</sup><sup>∗\*</sup>$``$Curiously enough, this type of reasoning closely resembles the notions of reduction and normalization for $`\lambda `$-calculus; it might be interesting to explore the analogy. that if two different symbols $`(r,u)`$ and $`(r^{},u^{})`$ are both beyond $`(s,f)`$, then there is a symbol $`(q,v)`$ which is beyond both $`(r,u)`$ and $`(r^{},u^{})`$. In this case we have morphisms $`a`$ and $`a^{}`$ serving as intermediaries between $`(s,f)`$ and $`(r,u)`$ or $`(r^{},u^{})`$ respectively. We would like to complete $`a`$ and $`a^{}`$ to a commutative square. For this we would hope to use condition (c), which requires one of the morphisms to be in $`\mathrm{\Sigma }`$. If we had condition (e) then this would be OK; in general an additional step is necessary.
Say that $`(r,u)`$ is under $`(s,f)`$ if there is a morphism $`a`$ intermediary from $`(s,f)`$ to $`(r,u)`$ such that $`a\mathrm{\Sigma }`$. Note that “under” implies “beyond” but not necessarily vice-versa. The main observation is the following lemma, which gives a sort of weak replacement for the $`3`$ for $`2`$ property, and its corollaries.
###### Lemma 4.1.
Suppose $`s\mathrm{\Sigma }`$ and $`a`$ is a morphism composable with $`s`$, such that $`r:=as`$ is in $`\mathrm{\Sigma }`$. Then there exists a morphism $`b`$ such that $`ba\mathrm{\Sigma }`$.
Proof: Consider the diagram
$$\stackrel{r}{}\stackrel{s}{}.$$
It is a right-fraction symbol because $`r\mathrm{\Sigma }`$. By condition (c) it can be transformed into a left-fraction symbol: there exist morphisms $`x`$ and $`t`$ with $`t\mathrm{\Sigma }`$ and $`xr=ts`$. We would like to factorize $`t`$ into a product $`ba`$, however we may need to go farther yet using condition (d). Our morphism $`a`$ goes from the target of $`s`$ to the target of $`r`$, and we have
$$xas=xr=ts.$$
In particular, we have two morphisms $`xa`$ and $`t`$ with the same source and target, equalized on the right by $`s\mathrm{\Sigma }`$. By condition (d) there is a morphism $`c\mathrm{\Sigma }`$ such that $`cxa=ct`$. Recall that $`t\mathrm{\Sigma }`$ so $`ct\mathrm{\Sigma }`$, and we can set $`b=cx`$ to obtain the lemma. $`\mathrm{}`$
In the proof files, the argument of Lemma 4.1 is integrated into the proof of Lemma `exists_lf_under` as in the following corollary.
###### Corollary 4.2.
Suppose $`(r,u)`$ is beyond $`(s,f)`$. Then there is another left fraction symbol $`(t,v)`$ such that $`(t,v)`$ is beyond $`(r,u)`$ and under $`(s,f)`$.
Proof: (see `exists_lf_under` in the proof files). Let $`a`$ be the intermediary morphism going from $`(s,f)`$ to $`(r,u)`$. Recall that $`r`$ and $`s`$ are in $`\mathrm{\Sigma }`$, and $`r=as`$. The lemma says there is another morphism $`b`$ such that $`ba\mathrm{\Sigma }`$. Let $`t=br=(ba)s`$ and $`v=bu=(ba)f`$. $`\mathrm{}`$
###### Corollary 4.3.
If $`(r,u)`$ and $`(t,v)`$ are both beyond $`(s,f)`$ then there is a symbol $`(q,w)`$ which is beyond $`(r,u)`$ and $`(t,v)`$.
Proof: By the previous corollary (and transitivity of “beyond” which is easy) we may assume that $`(r,u)`$ is under $`(s,f)`$. Then (and this part is Lemma `exists_lf_further` in the proof files) applying condition (c) to the intermediate morphisms we obtain intermediate morphisms going from $`(r,u)`$ and $`(t,v)`$ to a single $`(q,w)`$. $`\mathrm{}`$
Transitivity of the relation follows easily from Corollary 4.3. See `lf_equiv_trans` in the proof files.
A well-thought out direct argument for transitivity (which doesn’t occupy too much space) is given in Borceux Proposition 5.2.4. The essential information is reduced to a single diagram (Diagram 5.4, page 185) containing $`9`$ objects and $`13`$ arrows.
For the definition of the composition, Borceux writes (p. 185):
“…Moreover this definition is independent of the choices of $`f`$, $`s`$, $`g`$, $`t`$, $`h`$, $`r`$. This is lengthy but straightforward: the arguments are analogous to those for proving the transitivity of the equivalence relation defined on the arrows. We leave those details to the reader as well as the checking of the category axioms ….”
This analysis is basically sound: once one has gotten over the hurdle discussed above, which first shows up at the proof of transitivity, the remainder of the argument necessary for checking well-definedness of the composition, associativity and identity axioms and so forth, presents no further difficulties. Nonetheless, it might be the case that the simplified presentation of the proof of the transitivity of the equivalence relation, could have as a consequence that checking the facts about the composition law becomes more involved (we needed to use techniques similar to those for transitivity, in the proof of well-definedness of the composition for example). Similarly in and , the composition representative is constructed but well-definedness and associativity of the composition are not verified in detail.
For completeness, we describe here some of the main points. Suppose we are given two left-fraction symbols which are composable:
$$x\stackrel{f}{}y\stackrel{t}{}z\stackrel{g}{}u\stackrel{r}{}v.$$
Then the middle arrows give a right-fraction symbol which we can fill in to a square with a left-fraction symbol going in the other direction:
(4.6) $`\begin{array}{ccccc}& & z& & \\ & & & & \\ y& & & & u\\ & & & & \\ & & z^{}& & \end{array}`$
which in turn fits into the previous collection to yield a composite left-fraction symbol $`(g^{}f,t^{}r)`$:
$$x\stackrel{f}{}y\stackrel{g^{}}{}z^{}\stackrel{t^{}}{}u\stackrel{r}{}v.$$
In order to define the composition, we make a choice of fill-in square and set the composition equal to the composite symbol $`(g^{}f,t^{}r)`$. This composition rule is not associative, nor does it satisfy the left and right identity relations. On the other hand, modulo the equivalence relation established above, the composition will become associative and unitary.
The first and main step is to show that the composition is well-defined modulo the equivalence relation. This has two parts: first that if we make two different choices of fill-in square then the resulting composite symbols are equivalent; and secondly if we choose different representatives for the symbols which are being composed, then the composites are equivalent. These proofs make use of the same kind of arguments as we have described above, invoking things like Corollary 4.2 when necessary. The reader can by now imagine why no authors have attempted to write down the full text of these proofs in a forum destined for human readers. Those who are interested may refer directly to the proof files.
Once the well-definedness is established, the associativity is significantly easier at least on a conceptual level. It suffices to look at the following diagram:
(4.14) $`\begin{array}{ccccccccccccc}& & & & & & & & & & & & \\ & & & & & & & & & & & & \\ & & & & 1& & & & 2& & & & \\ & & & & & & & & & & & & \\ & & & & & & 3& & & & & & \\ & & & & & & & & & & & & \\ & & & & & & & & & & & & \end{array}`$
where, along the top, are the three left-fraction symbols we want to compose. Choose the top two fill-in squares denoted $`1`$ and $`2`$ first, which gives the middle row of arrows; then choose the bottom fill-in square $`3`$. Now the two different associated products may be obtained as follows (here we use the invariance under choice of fill-in square):
—one is obtained by using square $`1`$ to compose the first two symbols; then the composite rectangle $`2+3`$ is a fill-in square for multiplying this first composite with the rightmost symbol;
—the other is obtained by using square $`2`$ to compose the second two symbols; and the composite rectangle $`1+3`$ is a fill-in square for multiplying the leftmost symbol with this first composite.
Both methods give as result the left-fraction symbol obtained using the composites along the bottom edges of the big diagram. Thus, with the choices made as described above, the composition becomes associative “on the nose”; and because of the invariance of choices up to equivalence, we get that composition is associative up to equivalence when we make arbitrary choices for the fill-in squares.
The left and right unit conditions are proved similarly.
We obtain a category of “left fractions”. Defining the functor from our original category into the left fraction category, and proving that the images of elements of the localizing system are invertible, involve again some lemmas of a similar nature, whose proofs basically consist of setting up the appropriate diagrams and using good choices for the fill-in squares to define the compositions in question. For all of these things, we are in agreement with all of the authors found so far, that it isn’t worthwhile to write a mathematical text for these proofs. The proofs may be found directly in the proof files attached to .
The construction of the localization by left fractions can be considered as the statement of a nontrivial theorem about any localization (for example, about the general localization constructed previously).
###### Theorem 4.4.
Suppose $`𝒞`$ is a category and $`\mathrm{\Sigma }`$ is a multiplicative system satisfying the left-fraction conditions (a)–(d) above. Let $`𝒞[\mathrm{\Sigma }^1]`$ be a localization. Then the morphisms of $`𝒞[\mathrm{\Sigma }^1]`$ have the following description. Every morphism can be written as a composition $`t^1u`$ where $`u`$ comes from $`𝒞`$ and $`t`$ comes from $`\mathrm{\Sigma }`$ (and their targets coincide). If, for two such pairs $`t^1u=r^1v`$, then there exist morphisms $`a`$ and $`b`$ in $`𝒞`$ such that: $`a`$ is composable with $`u`$ and $`t`$; $`b`$ is composable with $`v`$ and $`r`$; $`au=bv`$; and $`at=br`$ and this is in $`\mathrm{\Sigma }`$.
The proof is that this description holds by definition for the left-fraction localization we are discussing in the present section. Then the universal properties show that any two localizations are isomorphic, so the same description holds in any other localization. This theorem is treated in the file `gzloc.v` (it is only there that we treat the fact that different localizations are isomorphic). It might be interesting to try to prove this description directly for the general construction of the localization. The left-fractions conditions imply fairly directly that morphisms in the localization can be written as simple products. However, to verify the statement about the equivalence relation seems difficult.
As a conclusion to this section, it is interesting to note that the mathematics behind the fraction construction is not one hundred percent straightforward, as was the case for the mathematics behind the general construction. On the other hand, it is commonly believed that the “calculus of fractions” construction is much more concrete and easy to understand. A possible reason for this is that mathematicians are very attached to considering the “size” of the mathematical objects which they manipulate, rather than the size of the associated mathematical theories. Since the arrows in the fraction construction are paths of length two whereas the arrows of the general construction are paths of arbitrary length, people prefer to think about the fraction construction (for example D. Pronk generalized the fraction construction to the case of $`2`$-categories but didn’t mention generalizing the general construction). This tendancy is similar to the constructionist or intuitionist philosophy: even while admitting reasoning based on less constructive arguments, mathematicians of all philosophies gravitate towards smaller and more constructive objects when they are available.
## 5. Formalizing the left-fraction construction
The formalization is contained in the file `left_fractions.v`, where `Left_Fractions` is the first module treating all of the essential constructions and properties. It starts with what is by now a fairly standard kind of definition, `lf_symbol f t` is an object (`Arrow.like`, in fact) containing the pair $`(f,t)`$ and corresponding to the left-fraction symbols used in the informal discussion above. The construction `lf_choice a s r g` represents a choice of fill-in left-fraction symbol creating a commutative square whose upper sides are the right-fraction symbol $`(r,g)`$.
A left-fraction symbol has an additional object besides its `source` and `target`, which we call `lf_vertex u`. This is the common target of the two morphisms involved. The operation `lf_extend a s p u` corresponds to composing both arrows of the left-fraction symbol $`u`$, with a morphism $`p`$ whose source is `source p = lf_vertex u`.
The `lf_extend` enters into the definitions of the notions `lf_beyond` and `lf_under` as defined in the previous section. In turn we define `lf_equiv` as existence of a common symbol which is `lf_beyond` the two in question. Then comes the main part of the proof which is Lemma `lf_equiv_trans`. This proof is done as described in the previous section (the division into sublemmas is slightly different from what is done informally above; we have referred above to the corresponding places in the proof files).
The main thing I would like to talk about in this section is the method we use for representing situations which, in informal argument, would be represented by a commutative diagram drawn in the text. Consider for example the definition
```
Definition fills_in a s u v w :=
has_left_fractions a s &
is_lf_symbol a s u & is_lf_symbol a s v & is_lf_symbol a s w
& source u = target v & source w = lf_vertex v &
target w = lf_vertex u &
comp a (lf_forward w) (lf_backward v) =
comp a (lf_backward w) (lf_forward u).
```
This represents the diagram 4.6 we have drawn in the previous section for defining composition. The variables `u`, `v` and `w` are the three left-fraction symbols occuring in the diagram (the first two on the upper row and the last one completing the bottom).
In a similar way, the definition `assoc_board a s u v w x y z` represents the diagram 4.14 we have drawn above for the associativity of composition. Another important pair of diagrams are `lf_lean_to a s e f g h i j` and `closes_lf_lean_to a s e f g h i j k l`. These correspond to diagrams, vaguely shaped like “lean-to’s”, which we haven’t drawn above (due mostly to my lack of -nique), and which enter into the proof that the composition is well-defined up to equivalence.
The idea in all of these cases is to make a definition involving all of the objects occurring in the diagram, which corresponds to commutativity of the diagram plus all of the other basic information it is supposed to satisfy (for example saying what the elements are, and that the sources and targets match up). In the cases `fills_in` and `assoc_board` we have chosen to represent the elements as being the left-fraction symbols (i.e. they are pairs of arrows) whereas in the diagrams `lf_lean_to` and `closes_lf_lean_to` the elements are morphisms of `a`. In both cases the first variable `a` is the category and the second variable `s` is the multiplicative system we are considering.
Given this way of manipulating diagrams, we can then state the main steps in the proof. For example, the main step of the well-definedness of composition is
```
Lemma lf_lean_to_closure : forall a s e f g h i j,
lf_lean_to a s e f g h i j ->
(exists k, exists l,(closes_lf_lean_to a s e f g h i j k l)).
```
This lemma is then used in Lemma `weak_rep_lf_equiv` (this fact doesn’t show up in the text as recopied in the preprint because it is inside the proof—which shows a limitation to the idea of just copying definitions and lemma statements as a simplified presentation).
Often the diagram definition will occur as a hypothesis of a lemma. This is particularly true of the hypothesis `assoc_board a s u v w x y z` which occurs as a hypothesis in a few different intermediate lemmas before the main statement of Lemma `make_comp_assoc_board`. The definitions `ffb_symbol` and `fbb_symbol` correspond to the two operations of juxtaposing two fill-in squares to get a rectangle which composes to a new fill-in square giving the outer compositions in the associativity statement (in the previous section this was where we looked at rectangles denoted $`1+3`$ or $`2+3`$).
As a general matter, writing mathematics for the computer requires that we go to a notation which is completely precise. Different strategies might be used for trying to keep a lid on the length of such notation. It is interesting to note that in situations such as the present one, precise definitions such as `ffb_symbol a s y z` replace vague statements such as “juxtapose the squares denoted $`1`$ and $`3`$ in the above diagram”. One major problem in both cases is the problem of refering to pieces of the diagram in question. The fact that we have included the various pieces as variables in our diagram definition `assoc_board a s u v w x y z` means that we can give small variable letters to each piece. Thus the `y` and `z` in `ffb_symbol a s y z` refer to places in `assoc_board`. In a mathematical text this referentiation operation becomes cumbersome when we start to manipulate large numbers of objects: we are led to circumlocutions like “the arrow on the upper left” and generally get lost in the meanders of referencing conventions of natural language (which are not totally precise nor sufficiently powerful). Another tempting solution would be to create a distinct mathematical object for the whole diagram. This is pretty much what we have done with the notion of `lf_symbol` for example. The drawback is that we then need long names for the component pieces of the diagrams. In the case of `lf_symbol` these were the functions `lf_forward` and `lf_backward`. This approach was called for in the case of `lf_symbol` because of the frequency and diversity of manipulations we have to do with these objects. On the other hand, with big diagrams which occur basically only once or a few times, it seems better to avoid allocating specific long names to their component pieces, and let the components be variables in a propositional representation of the diagram as a functional property. The observations in this paragraph are not intended as strict edicts but rather as ideas for one possible way to approach the language problems which are posed by computer formulation.
Rather than going on in detail about the remainder of the construction (the module `Left_Fraction_Category` is where we use the `Associating_Quotient` module to actually construct the category of left fractions; then we need to establish its universal properties and so forth), we close this section with an observation about proof technique. The proofs of the main lemmas referred to above are often rather long, because they involve manipulations of large amounts of information. Since steps such as rewriting tend to produce residual goals, one arrives at a situation where there is an impossibly large number of residual goals to treat at the end of a proof. Furthermore the treatment of these goals tends to be highly repetitive. In this situation it is essential to maintain a certain level of discipline in the following sense: when one comes upon a rewriting situation which generates an additional goal, one must go back to the start of the proof and add in that goal as an `Assert` statement. This can be done without changing the numbering of hypotheses (which would be painful to correct at each occurence of this phenomenon) by considering the `Assert` statements as “sublemmas” in the proof and naming them as such. Thus, rather than writing
```
Assert (...)
```
it is better to write
```
Assert (lemA : ...).
```
The proof itself becomes a location where there are many sublemmas. One can even recopy the sublemma texts with their proofs, from one proof to another (when the proof contexts are going to be similar in both cases). Once the required sublemmas are there, a combined rewriting tactic (such as our abbreviations `rw` or `wr`) which does a rewrite and then tries `assumption` and `trivial` on the subgoals, gives a proof where the number of auxiliary subgoals is reduced significantly enough to be only “very annoying” rather than “impossibly painful”.
## 6. Further formalizations
After the file on the left fraction construction, we include a few more files in the present development. In `gzloc.v`, we start by going back to some general considerations about localization. Lemmas whose names finish with `_recall` are statements meant to recall definitions from earlier files, so as to make reading of this part a little bit more self-contained. The first main corollary of the section is that two different localizations are isomorphic. Recall that `are_finverse a b` means that two functors are inverse (which is to say that their compositions are equal to the identity—hence they establish an isomorphism between their sources and targets). Thus the lemma `are_finverse_dotted_choice` says that any two localization functors are isomorphic (refer to `source_dotted_choice` and `target_dotted_choice` as well as `fcompose_dotted_choice` to complete this statement).
The next step is to investigate the relationship between all of our various definitions, and the opposite category and opposite functor structures. The main point here is to define the calculus of right fractions, by conjugating left fractions with “opposite”. This is of course completely straightforward, but many statements require lengthy proofs, which is certainly evidence that our overall setup is not really optimal.
In the part of the `gzloc.v` file starting with the definitions `lf_vee`, `lf_vee_image` and `lf_vee_equivalent`, we treat Theorem 4.4. As pointed out above, the proof is based on the fact that our left-fraction construction of a localization is isomorphic to any other localization due to the universal property. The statement of Theorem 4.4 is contained in Lemmas `left_fraction_description_for_loc` (for the case of an arbitrary localization) and `left_fraction_description_gz_proj` (for the case of the original general construction). Dualizing (by applying the opposite construction discussed in the previous paragraph) we obtain the corresponding results for right fractions.
The last file of the localization discussion is `lfcx.v`. Here we construct a little counterexample to one of the technical points encountered in the left-fractions construction. This part of the formalisation takes us back to one of the basic points of our approach to types and set-theory, namely that we integrate the inductive creation of types in Coq into our axiomatization of set theory. This allows us to manipulate small finite sets by creating them as inductive objects. In this way we can do a relatively large amount of case analysis by defining recursive tactics using the `Ltac` tactic language. In this way many of the proofs of properties of our constructed objects are very short lists of tactics (which take some time for the computer to digest). This could perhaps be thought of as a very very baby version of Gonthier-Werner’s techniques which they have applied to the 4-color theorem. The conclusion of the file is existence of a category `a` and localizing system `s`, satisfying left fractions, but with two left-fraction symbols `u` and `v` with `lf_beyond a s u v` (thus `u` and `v` project to equivalent arrows in the localization), however there is no symbol `w` such that `lf_under a s u w` and `lf_under a s v w`. This shows that we need to use the condition `lf_beyond` rather than `lf_under` (see the discussion of this point above).
The last file of the development is `infinite.v`. This is a complete digression from the subject: here we prove the fundamental properties of cardinal arithmetic for infinite cardinals, namely that the union or product of two infinite cardinals has the maximum of the two cardinalities. By Russell’s paradox on the other hand, the powerset of an infinite cardinal is strictly bigger. Along the way we do a certain number of basic properties of finite and infinite sets, cardinals, and ordinals. It is beyond the scope of the present preprint to go into further detail about the strategy of proof; and (contrarily to the case of localization of categories as we have seen) this is a subject in which there is no lack of different treatments in the literature—which we have not tried at all to index in the references.
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# Quantitative recurrence and large deviations for Teichmuller geodesic flow
## 1. Introduction
Let $`\mathrm{\Sigma }_g`$ be a compact surface of genus $`g2`$. Let $`Q_g`$ be the moduli space of unit-area holomorphic quadratic differentials on $`\mathrm{\Sigma }_g`$. That is, a point $`qQ_g`$ is a equivalence class of pairs $`(M,\omega )`$, where $`M`$ is a genus $`g`$ Riemann surface, and $`\omega `$ is a holomorphic quadratic differential on $`M`$, i.e., a tensor with the form $`f(z)dz^2`$ in local coordinates, such that $`_M|\omega |=1`$. Two pairs $`(M_1,\omega _1)`$ and $`(M_2,\omega _2)`$ are equivalent if there is a biholomorphism $`f:M_1M_2`$ such that $`f_{}\omega _1=\omega _2`$.
Given a pair $`qQ_g`$, one obtains (via integration of the form) an atlas of charts to $`^2`$, with transition maps of the form $`z\pm z+c`$. Similarly, given such an atlas of charts, one obtains a holomorphic quadratic differential by pulling back the form $`dz^2`$ on $``$.
These charts allow us to define a $`SL(2,)`$ action on $`Q_g`$ (and $`\stackrel{~}{Q_g}`$) given by linear post-composition with charts.
$`Q_g`$ is naturally the unit cotangent bundle to $`_g`$, the moduli space of Riemann surfaces. The fiber over each point $`M_g`$ is the vector space of holomorphic quadratic differentials on $`M`$.
The space $`Q_g`$ is naturally stratified by integer partitions $`\beta `$ of $`4g4`$. Strata are not always connected, however, they have at most finitely many components , and are invariant under $`SL(2,)`$. For the rest of our paper we work with one of these connected components, call it $`Q`$.
Without loss of generality, we will study strata of squares of abelian differentials (holomorphic 1-forms). Otherwise, we pass to a double cover. Since our results hold for all $`g2`$, we simply consider the (higher-genus) stratum defined in this way. Each stratum $`Q`$, while non-compact, is endowed with a continuous, ergodic, $`SL(2,)`$-invariant probability measure $`\mu _Q`$.
More details about the basics of moduli spaces and quadratic differentials can be found in, e.g., .
We are particularly interested in the action of the standard subgroups $`K`$ (maximal compact) and $`A`$ (diagonal matrices) of $`SL(2,)`$. These are both one-parameter subgroups, and can be described as follows:
$$K=\{r_\theta =\left(\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta \\ \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right):0\theta <2\pi \}$$
$$A=\{g_t=\left(\begin{array}{cc}e^t& 0\\ 0& e^t\end{array}\right):t\}.$$
We will often be interested also in the action of the semigroup $`A^+=\{g_t:t0\}`$.
The action of $`K`$, known as the *circle flow*, preserves the underlying holomorphic structure, so it acts as identity when projected to $`_g`$. The action of $`A`$ is known as *Teichmuller geodesic flow*, since the projection of an $`A`$-orbit yields a geodesic in the Teichmuller metric on $`_g`$ (and in fact, all Teichmuller geodesics arise this way).
Masur and Veech independently showed the Teichmuller geodesic flow is ergodic with respect to $`\mu _Q`$ and even mixing. Veech showed that it was ‘measurably Anosov’, and more recently Forni has obtained explicit formulas regarding the hyperbolic behavior. His results, which imply that as long as trajectories remain in a compact set their hyperbolicity can be controlled, provide much of the motivation for our research.
Our main results are concerned with the following scenario. Fix $`qQ`$, and consider the ‘circle’ $`Kq=\{r_\theta q:0\theta <2\pi \}`$. This set is endowed with a natural probability measure $`\nu `$, coming from the Haar measure on $`KS^1`$. What is the recurrence behavior of a trajectory $`\{g_tr_\theta q\}`$ when $`\theta `$ is chosen at random according to $`\nu `$?
This type of situation was first considered by Kerckhoff-Masur-Smillie , in order to understand dynamics of billiards in rational angle Euclidean polygons. They developed a dictionary between the dynamics of the straight-line flow on the singular Euclidean surface determined by a quadratic differential $`q`$ and the recurrence behavior of the geodesic trajectory $`g_tq`$ in the stratum $`Q`$. The $`K`$-action here corresponds to changing the direction of the straight line flow without changing the underlying surface. Thus, making a statement about a fixed $`q`$ and almost all directions allowed one to say something about straight line flow in almost all directions for a given flat surface, in analogy with Weyl’s equidstribution theorem on the torus.
To apply these results to billiards, one follows the unfolding procedure of Zemljakov and Katok , which translate questions about dynamics of billiards in rational angle polygons to that of flows on an associated singular Euclidean surface. However, the set of surfaces arising from billiards is of measure 0 in every stratum, so making statements about almost every point $`q`$ does not suffice. To make statements about billiards, we need statements that hold for all $`qQ`$.
The main result in is that for all $`qQ`$, and almost every $`\theta [0,2\pi )`$, the geodesic trajectory $`\{g_tr_\theta q\}_{t0}`$ is recurrent in $`Q`$. As a corollary, one obtains that the directional flow for Euclidean polygonal billiards is uniquely ergodic for almost every direction. Further results in this direction include estimates on the Hausdorff dimension of divergent and bounded trajectories (), and further relations between dynamics of the straight line flow and recurrence of the associated trajectory .
Our main results concern finer recurrence behavior of geodesic trajectories. In particular, for any fixed $`qQ`$, we estimate the measure of the set of angles such that the associated $`A^+`$ orbit $`\{g_tr_\theta q\}_{t0}`$ ‘behaves poorly’ for a length of time $`T`$. For us, poor behavior means that the trajectory is spending a lot of time in a neighborhood of the cusp of $`Q`$. Masur proved a statistical result, known as a *logarithm law* in this situation, in analogy with earlier results on symmetric spaces due to Sullivan and Kleinbock-Margulis .
We will construct a proper (that is, unbounded off compact sets) continuous function $`V:Q^+`$, and consider the recurrence behavior of $`\{g_tr_\theta q\}_{t0}`$ to the family of compact sets $`C_l:=\{q:V(q)l\}`$, which form an exhaustion of $`Q`$ as $`l`$ varies. Our main results can be summarized as follows:
###### Theorem 1.1.
Fix notation as above. Then
1. For all $`l`$ sufficiently large and all $`qC_l`$, there are positive constants $`c_1=c_1(l,q),c_2(l)`$, with
$$\nu \{\theta :g_tr_\theta qC_l,0tT\}c_1e^{c_2T}$$
for all T sufficiently large. That is, the probability that a random geodesic trajectory has not visited $`C_l`$ by time $`T`$ decays exponentially in $`T`$.
2. For all $`l,S,T`$ sufficiently large and all $`qQ`$, there are positive constants $`c_3=c_3(S,l,q),c_4=c_4(l)`$, with
$$\nu \{\theta :g_tr_\theta qC_l,StS+T\}c_3e^{c_4T}.$$
That is, the probability that a random geodesic trajectory does not enter $`C_l`$ in the interval $`[S,S+T]`$ decays exponentially in $`T`$.
3. Let $`qQ`$. For any $`0<\lambda <1`$, there is a $`l0`$, and $`0<\gamma <1`$, such that for all $`T`$ sufficiently large (depending on all the above constants)
$$\nu \{\theta :\frac{1}{T}|\{0tT:g_tr_\theta qC_l\}|>\lambda \}\gamma ^T.$$
Result (3) above may be thought of as a large deviations result for the Teichmuller flow. While ergodicity guarantees that $`\frac{1}{T}\left|\{0tT:g_tqC_l\}\right|\mu _Q(C_l)`$ for $`\mu _Q`$-almost every $`qQ`$, our result gives explicit information about the likelihood of bad trajectories. Notice, however, this is *not* a traditional large deviations result, which estimates the probability of a deviation of any $`ϵ>0`$ from the ergodic average. Other interesting results concerning deviations are due to Bufetov , who proved a central limit theorem for this flow.
In , Forni related this type of fine recurrence behavior for the geodesic trajectory $`\{g_tq\}_{t0}`$ to deviation of ergodic averages for the straight line flow on the flat surface associated to $`q`$. He proved that for almost every $`qQ`$, the associated flow has the same deviation behavior. However, his result gives no information on billiards, as the set of quadratic differentials arising from billiards are a set of measure zero. Since our results provide fine recurrence information about the geodesic trajectory $`\{g_tr_\theta q\}_{t0}`$ for *all* $`qQ`$, and almost all $`\theta `$, we conjecture the following:
###### Conjecture.
For all rational-angle Euclidean polygons, the deviation of ergodic averages for the billiard flow is the same for almost all directions, and depends only on the $`SL(2,)`$ orbit of the associated quadratic differential.
It was also shown in that as long as a geodesic trajectory stays within a compact set, the rate of expansion/contraction in the tangent space along the trajectory is bounded away from $`1`$. Thus our results can be used to obtain explicit estimates on the hyperbolicity of the flow along specific trajectories. Avila-Gouezel-Yoccoz used exponential return estimates to a different family of compact sets to prove the exponential rate of mixing for the Teichmuller geodesic flow, which was the original motivation for this research.
Other quantitative recurrence results for dynamics on Teichmuller spaces were obtained by Minsky-Weiss for the case of Teichmuller horocycle flow.
We also have a collection of results for a certain class of random walks, defined as follows: Fix $`\tau >0`$. Given that we are at a point $`qQ`$, the next point in our trajectory will be chosen at random according to Haar measure on $`S^1`$ from the ‘circle’ of radius $`\tau ,\{g_\tau r_\theta q:0\theta <2\pi \}`$. Note that $`\mu _Q`$ is a stationary measure for this walk, since it is $`SL(2,)`$-invariant. Let $`\{X_n\}_{n=0}^{\mathrm{}}`$ denote the random walk generated this way.
Remark: By hyperbolic geometry, one can see that trajectories of this walk closely approximate geodesics. Thus, understanding recurrence properties of the walk gives one insight into properties of the flow.
We define $`P_q(E):=\text{Prob}(E|X_0=q)`$ for any event $`E`$ defined on the trajectory $`\{X_n\}_{n=0}^{\mathrm{}}`$ starting at $`X_0=q`$. Also, for any measurable $`AQ`$, we define $`P^n(q,A):=P_q(X_nA)`$. That is, $`P^n(q,.)`$ is the probability distribution of $`X_n`$ given $`X_0=q`$.
###### Theorem 1.2.
For all $`\tau `$ sufficiently large, we have:
1. There is an $`l(\tau )`$ such that for all $`l>l(\tau )`$, and all $`qC_l`$, there are constants $`c_5=c_5(q,l,\tau ),c_6=c_6(l,\tau )`$, such that
$$P_q(X_iC_l:1in)c_5e^{c_6n}.$$
2. For all compact $`CQ`$, $`ϵ>0`$, there is a $`C_lC`$ such that $`qC`$, and for all $`m0`$,
$$P_q(X_mC_l)>1ϵ.$$
3. For all $`ϵ>0`$, there is a $`l=l(ϵ)>0`$ such that $`qQ`$, there is an $`M(q)`$ such that for all $`m>M(q)`$,
$$P_q(X_mC_l)>1ϵ.$$
The rest of the paper is organized as follows: In the next section, we give more detailed statements of our main results, and construct the function $`V`$. We also give a version of our results for general $`SL(2,)`$-actions. In section 3, we prove our results for random walks. In section 4, we collect some technical lemmas about change of polar coordinates on hyperbolic space. In section 5, we construct the required lemmas from the theory of large deviations. In section 6, we prove our main theorems for the flow.
## 2. Statement of Results
The precise statements for the flow regard a *family* of proper functions $`V_\delta `$, $`0<\delta <1`$, and the compact sets $`C_{\delta ,l}:=\{q:V_\delta (q)l\}`$. We have:
### 2.1. Flow results
###### Theorem 2.1.
For every $`1>\delta >0`$ there is a proper (i.e., unbounded off compact sets), smooth, $`K`$-invariant function $`V_\delta :Q^+`$ and positive constants $`t_0=t_0(\delta ),l_0=l_0(\delta ),a=a(\delta )`$ such that for all $`ll_0`$, there are $`1>\delta ^{}=\delta ^{}(l,\delta )>\delta `$, with $`\delta ^{}`$ decreasing as a function of $`l`$, so that for all $`qC_{\delta ,l}`$
$$\nu \{\theta :g_tr_\theta qC_{\delta ,l},0tT\}a\frac{V_\delta (q)}{l}e^{(1\delta ^{})T},$$
for all $`T>t_0`$.
This is the precise version of part (1) of Theorem 1.1. Note that the outside term essentially depends only on $`V_\delta (q)`$.
###### Theorem 2.2.
Let $`qQ`$. For every $`ϵ>0`$, $`1>\delta >0`$, there are positive constants $`S_0,T_1,l_1`$ depending on $`\delta `$, $`1>\delta ^{\prime \prime }=\delta ^{\prime \prime }(l)>\delta `$, and $`\alpha =\alpha (q)`$ such that
$$\nu \{\theta :g_tr_\theta qC_{\delta ,l},StS+T\}\alpha e^{(1\delta ^{\prime \prime })T},$$
for all $`S>S_0`$, $`T>T_1`$, and $`l>l_1`$, with
$$\alpha 8(1+ϵ)\underset{\theta [0,2\pi )}{sup}ab\frac{V_\delta (g_Sr_\theta q)}{l},$$
where $`b`$ depends only on the curvature of $`^2`$.
Here, $`S_0`$ depends only on the choice of curvature for the hyperbolic plane $`^2`$, and $`T_1`$ is the maximum of $`t_0`$ from Theorem 2.1 and a $`T_2`$ depending only on curvature.
###### Theorem 2.3.
Let $`qQ`$. For any $`0<\lambda <1`$, and any $`0<\delta <1`$ there is a $`l0`$, and $`\gamma <1`$, such that
$$\nu \{\theta :\frac{1}{T}|\{0tT:g_tr_\theta qC_{\delta ,l}\}|>\lambda \}\gamma ^T,$$
for all $`T`$ sufficiently large.
This theorem uses a technical tool from the theory of large deviations, Proposition 5.1, which makes it difficult to track the dependence of $`\gamma `$ on $`l`$ and $`\lambda `$. Clearly, if an $`l`$ works for a fixed $`\lambda _0`$ and $`\delta `$, it works for all $`\lambda >\lambda _0`$, and any larger $`l`$ (with the same $`\delta `$) will work for $`\lambda _0`$. Similarly, the same $`l`$ will work for the same $`\lambda _0`$ and any smaller $`\delta `$.
In a personal communication, Forni posed the following question: for *every* $`q`$ in $`Q`$, and almost all $`\theta `$, does there exist a $`\zeta >0`$ such that the geodesic trajectory $`\{g_tr_\theta q\}_{t0}`$ spends at least a proportion $`\zeta `$ of its time in a fixed compact set? The following corollary answers in the affirmative.
###### Corollary 2.4.
Let $`qQ,\zeta >0`$. Fix $`ϵ>0`$. Let $`l`$ be such that Theorem 2.3 is satisfied with $`\lambda =1\zeta ,\delta =ϵ`$. Then for $`\nu `$-almost every $`\theta `$,
$$\underset{T\mathrm{}}{lim\; sup}\frac{1}{T}\left|\{0tT:g_tr_\theta qC_{\delta ,l}\}\right|\lambda .$$
The proof of this result follows from an application of the Borel-Cantelli lemma and Theorem 2.3.
### 2.2. Random walks on $`Q`$
Let $`\{X_n\}_{n=0}^{\mathrm{}}`$ be our trajectory. That is,
$$X_{n+1}=g_\tau r_{\theta _n}X_n,$$
where $`\{\theta _n\}_{n=0}^{\mathrm{}}`$ are independent and identically distributed (i.i.d.) according to the uniform distribution on $`[0,2\pi )`$ (equivalently, $`X_{n+1}=g_\tau k_nX_n`$, where $`\{k_n\}_{n=0}^{\mathrm{}}`$ are i.i.d. according to Haar measure on $`K`$). This implies $`\{X_n\}`$ is a Markov chain, with transition probability function $`P(x,A)=\nu \{\theta :g_\tau r_\theta qA\}`$.
###### Theorem 2.5.
Fix $`0<\delta <1`$. Let $`\{C_{\delta ,l}\}`$ be as in Theorem 2.1. Then there is a $`\tau _0=\tau _0(\delta )>0`$ such that for all $`\tau >\tau _0`$, there is a $`\stackrel{~}{l}_0=\stackrel{~}{l_0}(\delta ,\tau )`$ such that for all $`l>\stackrel{~}{l_0}`$, there is a $`\gamma =\gamma (l,\delta )<1`$ such that for all $`qC_{\delta ,l}`$, and for all $`n1`$,
$$P_q(X_jC_{\delta ,l}:0jn)\frac{V_\delta (q)}{l}\gamma ^n.$$
This is completely analagous to Theorem 2.1, and indeed the proof of this result will be essential to the proof of Theorem 2.1.
For the rest of the section, fix $`1>\delta >0`$, and $`\tau >\tau _0(\delta )`$.
###### Theorem 2.6.
Let $`ϵ>0`$. Then:
1. For all compact $`CQ`$, there is a $`C_{\delta ,l}C`$ such that $`qC`$, and for all $`m0`$,
$$P_q(X_mC_{\delta ,l})>1ϵ.$$
2. There is a $`l>0`$ such that $`qQ`$, there is an $`M(q)`$ such that for all $`m>M(q)`$,
$$P_q(X_mC_{\delta ,l})>1ϵ.$$
Finally, we have:
###### Corollary 2.7.
Fix $`qQ`$, $`0<\lambda <1`$, and $`ϵ>0`$. Then there is a $`l=l(q,\lambda ,\delta ,ϵ)>0`$ so that, for all $`n`$ sufficiently large,
$$P_q\left(\frac{1}{n}|\{1in:X_iC_{\delta ,l}\}|>\lambda \right)>1ϵ$$
This is analogous to Corollary 2.4.
Remark: $`l`$ can be chosen uniformly as $`q`$ varies over a compact set.
### 2.3. $`SL(2,)`$-actions
Let $`SL(2,)`$ act continuously on a topological space $`X`$. Suppose there is a family of $`K`$-invariant functions $`V_\delta :X`$, $`0<\delta <1`$ satisfying the following properties:
* For all $`xX`$, consider the function $`V_{\delta ,x}:SL(2,)^+`$ defined by $`V_{\delta ,x}=V_\delta (gx)`$. Note that by $`K`$-invariance, we can view $`V_{\delta ,x}`$ as a function on $`^2=SO(2)`$`\`$`SL(2,)`$. We require the following : For all $`\sigma >1`$, there exists a $`\kappa >0`$ such that for any $`p^2`$ with $`d(p,i)<\kappa `$, and any $`hSL(2,)`$
(2.8)
$$\sigma ^1V_{\delta ,x}(h)V_{\delta ,x}(ph)\sigma V_{\delta ,x}(h).$$
We say such a function is *logsmooth*.
* For all $`1>\delta >0`$, there is a constant $`\stackrel{~}{c}=\stackrel{~}{c}(\delta )`$ such that for all sufficiently large $`\tau `$, there is a $`\stackrel{~}{b}=\stackrel{~}{b}(\tau ,\delta )`$ such that for all $`xX`$,
(2.9)
$$(A_\tau V_\delta )(x):=_0^{2\pi }V_\delta (g_\tau r_\theta x)𝑑\nu (\theta )\stackrel{~}{c}e^{(1\delta )\tau }V_\delta (x)+\stackrel{~}{b}.$$
All our results on recurrence to the sublevel sets $`C_{\delta ,l}`$ from subsections 2.1 and 2.2, while constructed for the space $`Q`$ hold for any space $`X`$ with the above properties.
The key lemma needed to prove our results will be the construction of a family of functions $`\{V_\delta \}`$ on $`Q`$ satisfying the above requirements.
###### Lemma 2.10.
There is a family of smooth, proper functions $`V_\delta :Q^+`$, $`0<\delta <1`$ satisfying equations (2.8) and (2.9).
We require the following technical lemma from (page 465, Lemma 7.5).
###### Lemma.
For all $`\delta >0`$ there are logsmooth functions $`V_0,\mathrm{},V_n:Q^+`$ such that $`V_0`$ is proper, and for every $`\tau >0`$, there are constants $`w=w(\tau ,\delta ),\stackrel{~}{b^{}}=\stackrel{~}{b^{}}(\tau ,\delta )`$, and $`\stackrel{~}{c^{}}=\stackrel{~}{c^{}}(\delta )`$, *independent* of $`\tau `$ such that for all $`0in`$, and $`qQ`$,
$$(A_\tau V_i)(q)\stackrel{~}{c^{}}e^{(1\delta )\tau }V_i(q)+w\underset{j=i+1}{\overset{n}{}}V_j(q)+\stackrel{~}{b^{}}.$$
Remark: In fact $`V_0(q)=\mathrm{max}(1,\frac{1}{l(q)^{1+\delta }})`$, where $`l(q)`$ denotes the length of the shortest saddle connection on $`q`$. Recall that a *saddle connection* is a geodesic (in the metric determined by $`q`$) connecting two zeroes of $`q`$.
Proof of Lemma 2.10:
Fix $`\delta ,\tau >0`$. Let $`\stackrel{~}{c}=2\stackrel{~}{c^{}}`$ $`\lambda _0=\frac{w}{\stackrel{~}{c^{}}}`$, $`\lambda _i=\left(\frac{\stackrel{~}{c^{}}}{w}+1\right)^{i1}`$, for $`0in`$. Note that
$$\underset{i=0}{\overset{j1}{}}\lambda _i\frac{\stackrel{~}{c}}{2w}\lambda _j.$$
Set $`\stackrel{~}{b}=\stackrel{~}{b^{}}_{i=0}^n\lambda _i`$. Let $`V_\delta (q)=_{i=0}^n\lambda _iV_i(q)`$. Then
$`\left(A_\tau V_\delta \right)(q)={\displaystyle \underset{i=0}{\overset{n}{}}}\lambda _i\left(A_\tau V_i\right)(q)`$ $``$ $`{\displaystyle \underset{i=0}{\overset{n}{}}}\lambda _i\left(\stackrel{~}{c^{}}e^{(1\delta )\tau }V_i(q)+w{\displaystyle \underset{j=i+1}{\overset{n}{}}}V_j(q)+\stackrel{~}{b^{}}\right)`$
$`=`$ $`{\displaystyle \underset{i=0}{\overset{n}{}}}\lambda _i\stackrel{~}{c^{}}e^{(1\delta )\tau }V_i(q)+w{\displaystyle \underset{j=1}{\overset{n}{}}}\left({\displaystyle \underset{i=1}{\overset{j1}{}}}\lambda _i\right)V_j(q)+b`$
$`=`$ $`\stackrel{~}{c^{}}e^{(1\delta )\tau }\lambda _0V_0(q)+{\displaystyle \underset{i=1}{\overset{n}{}}}V_j(q)\left(\lambda _j\stackrel{~}{c^{}}e^{(1\delta )\tau }+w{\displaystyle \underset{i=0}{\overset{j1}{}}}\lambda _i\right)+\stackrel{~}{b}`$
$`=`$ $`\stackrel{~}{c^{}}e^{(1\delta )\tau }\lambda _0V_0(q)+{\displaystyle \underset{j=1}{\overset{n}{}}}2\lambda _j\stackrel{~}{c^{}}e^{(1\delta )\tau }V_j(q)+\stackrel{~}{b}`$
$``$ $`\stackrel{~}{c}e^{(1\delta )\tau }V_\delta (q)+\stackrel{~}{b}.`$
Thus we have constructed a family of functions $`V_\delta `$ satisfying equations (2.8) and (2.9). That $`V_\delta `$ satisfies (2.8) follows from the logsmoothness of the $`V_i`$’s noted in (where they are called $`\alpha _i`$’s) on page 471, at the beginning of the proof of Proposition 7.2. Finally, $`V_\delta `$ is proper since $`V_0`$ is proper. ∎
## 3. General Markov Chain results
Our goal in this section is to recall some results from the theory of Markov Chains which will allow us to prove Theorems 2.5-2.6. We also give some applications to random walks on homogeneous spaces, as considered in . Many of these results can be found, in greater generality, in .
Before stating our main results, we recall some basic notation: If $`\{X_n\}`$ is a Markov chain on $`(S,𝒮)`$, $`S`$ the state space, and $`𝒮`$ the $`\sigma `$-algebra, then, for any event $`E`$ (an event $`E`$ is a set in the product $`\sigma `$-algebra) and any starting point $`xS`$ , $`P_x(E):=P(E|X_0=x)`$, i.e., it is the probability of the event $`E`$ occuring given that our starting point was $`x`$. Similarly, given a measurable subset $`CS`$, we write $`P^n(x,C)=P_x(X_nC)`$, and we define the first hitting time of $`C`$ by $`\tau _C:=inf\{n1:X_nC\}`$.
The main result of this section is:
###### Proposition 3.1.
Let $`S`$ be a non-compact topological space, and $`𝒮`$ its Borel $`\sigma `$-algebra. Let $`\{X_n\}_{n=0}^{\mathrm{}}`$ be a Markov chain on $`(S,𝒮)`$. Suppose there exists a smooth, proper function $`V:S^+`$ and constants $`0<c<1`$ and $`b>0`$ such that
$$(PV)(x):=E(V(X_1)|X_0=x)cV(x)+b,$$
for all $`xS`$. Then
1. For any $`l>0`$ and $`xC_l:=\{yS:V(y)l\}`$,
$$p_n(x):=P_x\left(\tau _{C_l}>n\right)\frac{V(x)}{l}\left(c+\frac{b}{l}\right)^n,$$
for all $`n0`$.
2. For all $`ϵ>0`$, and for all compact $`CS`$, there is an $`l`$ such that $`xC,m0`$,
$$P^m(x,C_l)>1ϵ.$$
3. For all $`ϵ>0`$ there is an $`l>0`$ such that $`xS`$, there is an $`M(x)>0`$ such that for $`m>M(x)`$,
$$P^m(x,C_l)>1ϵ.$$
4. For all $`xS`$, $`0<\lambda <1`$, and $`ϵ>0`$, there is a $`l=l(q,\lambda ,ϵ)>0`$ so that, for all $`n`$ sufficiently large,
$$P_x\left(\frac{1}{n}|\{1in:X_iC_l\}|>\lambda \right)>1ϵ.$$
If we fix $`\lambda `$ and $`ϵ`$, then $`l`$ can be chosen uniformly as $`x`$ varies over a compact set.
Proof of Theorems 2.5-2.6 and Corollary 2.7 : Combine Proposition 3.1 with Lemma 2.10. For Theorem 2.5, fix $`\delta >0`$ and let $`\tau _0`$ be such that $`c=\stackrel{~}{c}e^{(1\delta )\tau _0}<1`$, we obtain our result with $`\gamma =\left(c+\frac{\stackrel{~}{b}}{l}\right)`$ for $`l>\frac{\stackrel{~}{b}}{1c}`$.∎
Proof of Proposition 3.1:
* Proof of (1): Let $`B_n:=\{\tau _{C_l}>n\}`$. Then $`p_n(x)=P_x(B_n)`$. For $`n0`$,
$$lp_nE_x(V(X_n):B_n)=:D_n$$
since on $`B_n`$, $`X_nC_l`$, i.e. $`V(X_n)l`$.
Now, $`B_nB_{n1}`$, so
$`D_n`$ $``$ $`E_x(V(X_n):B_{n1})`$
$`=`$ $`E_x(E(V(X_n)|X_{n1}):B_{n1})`$
$`=`$ $`E_x((PV)(X_{n1}):B_{n1})`$
where we are using the Markov property in the 2nd line.
Now we can apply our condition
$$(PV)(X_{n1})cV(X_{n1})+b.$$
This yields
$$D_ncD_{n1}+bp_{n1},$$
and using the observation that $`p_{n1}\frac{D_{n1}}{l}`$, we obtain the recurrence relation
$$D_n\left(c+\frac{b}{l}\right)D_{n1}.$$
Iterating this, we obtain
$$D_nD_0\left(c+\frac{b}{l}\right)^n.$$
Since $`D_0=V(X_0)=V(x)`$, and $`p_n\frac{D_n}{l}`$, we obtain our result. Note that the result is only meaningful if $`(c+\frac{b}{l})<1`$, which is equivalent to setting $`l>\frac{b}{1c}`$.
* Proof of (2): We have $`(PV)(x)cV(X)+b`$ for $`c<1`$. Iterating this, we get that $`(P^mV)(x)c^mV(x)+b^{}`$, where $`b^{}`$ does not depend on $`m`$ or $`x`$. Set $`l=sup_{yC}\frac{V(y)+b^{}}{ϵ}`$. Then we have that
$$lP^m(x,C_l^c)E_x(V(X_m))c^mV(x)+b^{},$$
so we get that $`P^m(x,C_l^c)<ϵ`$ as desired. ∎
* Proof of (3): For the third property, we select l = $`2b^{}/ϵ`$, where $`b^{}`$ is as above. For $`m`$ sufficiently large, $`c^mV(x)b^{}`$, so by the argument in part (2), we can get our conclusion. ∎
* Proof of (4): Fix $`xS`$, $`0<\lambda <1`$, and $`ϵ>0`$. By part $`(3)`$, there is an $`l>0`$ so that for $`n`$ sufficiently large, $`P(X_nC_l)>1ϵ^{}`$, where $`0<ϵ^{}<\frac{2}{3}ϵ(1\lambda )`$. Set $`S_n=\frac{1}{n}_{i=1}^n\chi _{C_l}(X_i)`$, where $`\chi _{C_l}`$ is the indicator function of $`C_l`$. Then, for any $`\lambda <1`$,
$`E(S_n)`$ $``$ $`\lambda P(S_n\lambda )+P(S_n>\lambda )`$
$`=`$ $`\lambda +(1\lambda )P(S_n>\lambda ).`$
Thus, we have
$$P(S_n>\lambda )\frac{E(S_n)\lambda }{1\lambda }.$$
Now, for $`n`$ sufficiently large $`E(S_n)1\frac{3}{2}ϵ^{}`$, thus,
$$P(S_n>\lambda )\frac{1\frac{3}{2}ϵ^{}\lambda }{1\lambda }>1ϵ.$$
The fact that $`l`$ can be chosen uniformly as $`x`$ varies over a compact set follows from part (2).∎
These types of questions were considered for random walks on homogeneous spaces in by Eskin and Margulis. They constructed a function $`V`$ on their state space satisfying the conditions of Proposition 3.1, and used this to draw conclusions (2) and (3). Conclusions (1) and (4) appear to be new results for these walks.
Precisely, we have the following:
###### Theorem.
Let $`G`$ be a semisimple Lie group, and $`\mathrm{\Gamma }`$ a non-uniform lattice. Let $`\mu `$ be a probability measure on $`G`$ satisfying the conditions of Theorem 2.1 in . Consider the Markov chain $`\{X_n\}_{n=0}^{\mathrm{}}`$ defined on $`G/\mathrm{\Gamma }`$ by the measure $`\mu `$:
$$X_{n+1}=g_nX_n,$$
with $`\{g_n\}_{n=0}^{\mathrm{}}`$ an i.i.d. (with distribution $`\mu `$) sequence of elements of $`G`$. Then there is a function $`V:G/\mathrm{\Gamma }^+`$ satisfying the conditions of Proposition 3.1. Thus, conclusions (1)-(4) of the Proposition are satisfied.
## 4. Polar coordinates and shadowing
We require two lemmas about change of polar coordinates in the hyperbolic plane $`^2=SO(2)`$`\`$`SL(2,)`$. We fix two positive numbers $`t_1,t_2`$, and basepoints $`i`$ and $`z_0=i.g_{t_1}r_\theta `$ (these will correspond to our basepoint $`q`$ and an arbitrary $`q_0`$ in its $`SL(2,)`$-orbit, projected to $`SO(2)`$`\`$`Q`$). We let $`d(.,.)`$ denote distance in the hyperbolic plane.
Consider the circle of radius $`t_2`$ around $`z_0`$, defined by $`\{z_\varphi =i.g_{t_2}r_\varphi g_{t_1}r_\theta :0\varphi <2\pi \}`$. We say that $`t_2,\varphi `$ are the polar coordinates of $`z_\varphi `$ based at $`z_0`$. For each $`\varphi `$, we define $`D=D_{t_1,t_2}(\varphi )`$ and $`\mathrm{\Psi }=\mathrm{\Psi }_{t_1,t_2}(\varphi )`$ by $`z_\varphi =i.g_{D(\varphi )}r_{\theta +\mathrm{\Psi }(\varphi )}`$, i.e. $`D(\varphi ),\theta +\mathrm{\Psi }(\varphi )`$ are the polar coordinates of $`z_\varphi `$ based at $`i`$. Note that $`D,\mathrm{\Psi }`$ are *independent* of $`\theta `$.
Geometrically, $`D(\varphi )`$ is the distance and $`\mathrm{\Psi }(\varphi )`$ is the angle (measured clockwise from the the geodesic connecting $`i`$ to $`z_0`$) of the geodesic segment connecting $`i`$ to $`z_\varphi `$. Hyperbolic trigonometry (the laws of sines and cosines, appied to the triangle formed by the points $`i`$, $`z_0`$, and $`z_\varphi `$) yield:
(4.1)
$$\mathrm{cosh}D(\varphi )=\mathrm{cosh}t_1\mathrm{cosh}t_2+\mathrm{sinh}t_1\mathrm{sinh}t_2\mathrm{cos}\varphi ,$$
and
(4.2)
$$\mathrm{sin}\mathrm{\Psi }(\varphi )=\frac{\mathrm{sinh}t_2}{\mathrm{sinh}D(\varphi )}\mathrm{sin}\varphi .$$
If $`t_2>t_1`$, $`i`$ lies inside the circle of radius $`\tau `$ around $`z_0`$, and thus, the map $`\mathrm{\Psi }`$ is both one-to-one and onto. If $`t_1>t_2`$, the point $`i`$ is outside the cirlce, and $`\mathrm{\Psi }`$ is neither one-to-one or onto. In this case, the image is an interval, with boundary points such that the geodesic determined by those angles intersects the circle of radius $`t`$ tangentially. In the interior of the interval, each point $`\psi `$ has two preimages, call them $`\varphi _1,\varphi _2`$ one such that $`D(\varphi _1)t_1t_2`$, and one such that $`D(\varphi _2)t_1+t_2`$. For our applications, we will only be concerned with $`\varphi _2`$.
The key technical lemma is as follows:
###### Lemma 4.3.
Let $`A[0,2\pi )`$ be a measurable set. Then, for every $`ϵ>0`$, there are $`\tau _1,\tau _2`$, such that for all $`t_1>\tau _1,t_2>\tau _2`$ the neighborhood $`U=\mathrm{\Psi }_{t_1,t_2}([\pi /2,\pi /2])`$ of $`0`$ satisfies
$$\frac{\nu (U\mathrm{\Psi }(A))}{\nu (U)}4(1+ϵ)\nu (A).$$
Proof: For this estimate, we need to control the behavior of the derivative $`\mathrm{\Psi }^{}`$. More precisely, we need to control *ratios* $`\mathrm{\Psi }^{}(\varphi _1)/\mathrm{\Psi }^{}(\varphi _2)`$, with $`\varphi _1,\varphi _2\mathrm{\Psi }^1U`$, so we can compare $`\nu `$ and $`\mathrm{\Psi }_{}\nu `$, where $`\mathrm{\Psi }_{}\nu (E)=\nu (\mathrm{\Psi }^1E)`$.
We have the following claim:
###### Claim.
Let $`\eta >0`$. For $`t_1,t_2`$ sufficiently large,
$$\frac{e^{t_1}}{2}(1\eta )|\mathrm{\Psi }^{}(\varphi )|e^{t_1}(1+\eta ),$$
for $`\varphi [\pi /2,\pi /2]`$.
Proof of Claim:
Implicit differentiation of equations (4.1) and (4.2) yield:
(4.4)
$$D^{}(\varphi )\mathrm{sinh}D(\varphi )=\mathrm{sinh}t_1\mathrm{sinh}t_2\mathrm{sin}\varphi $$
and
(4.5)
$$\mathrm{\Psi }^{}(\varphi )\mathrm{cos}\mathrm{\Psi }(\varphi )=\mathrm{sinh}t_2\frac{\mathrm{cos}\varphi \mathrm{sinh}D(\varphi )+\mathrm{sin}^2\varphi \mathrm{coth}D(\varphi )\mathrm{sinh}t_1\mathrm{sinh}t_2}{\mathrm{sinh}^2D(\varphi )}.$$
Let $`\kappa >0`$. Let $`t_1,t_2`$ be large enough so that for all $`\varphi [\pi /2,\pi /2]`$,
1. $`D(\varphi )>t_1+t_2\kappa `$,
2. $`\mathrm{coth}(D(\varphi ))>1\kappa `$
3. $`\mathrm{cos}(\mathrm{\Psi }(\varphi ))>1\kappa `$
4. $`|\mathrm{sinh}t_1\mathrm{sinh}t_2\mathrm{sinh}D(\varphi )|\kappa `$
5. $`1\kappa <\frac{2\mathrm{sinh}t_1}{e^{t_i/2}},\frac{2\mathrm{cosh}t_1}{e^{t_i/2}}<1+\kappa `$ for $`i=1,2`$.
That we can achieve the above inequalities follows from hyperbolic geometry and the basic properties of $`\mathrm{sinh}`$ and $`\mathrm{cosh}`$.
Let $`\eta >0`$. Using the above inequalities and some basic algebra, we can select $`\tau _1,\tau _2`$ such that for all $`t_1>\tau _1,t_2>\tau _2`$, we have
$$(1\eta )\frac{\mathrm{\Psi }^{}(\varphi )}{\frac{e^{t_1}}{2}(2\mathrm{cos}\varphi +\mathrm{sin}^2\varphi )}(1+\eta ).$$
Remark: The expression $`\frac{e^{t_1}}{2}(2\mathrm{cos}\varphi +\mathrm{sin}^2\varphi )`$ is obtained by replacing the quantities in equation (4.5) with their approximations (1)-(5).
Now, let $`f(\varphi )=2\mathrm{cos}\varphi +\mathrm{sin}^2\varphi `$. We have $`1f(\varphi )2`$ for $`\varphi [\pi /2,\pi /2]`$. Thus,
$$\frac{e^{t_1}}{2}(1\eta )\mathrm{\Psi }^{}(\varphi )(1+\eta )e^{t_1},$$
completing the proof of the claim.
To complete the proof of the lemma, let $`ϵ>0`$ and $`\eta `$ be such that $`\frac{1+\eta }{1\eta }1+ϵ`$. By the claim, we know that the proportion of measure of any set in $`[\pi /2,\pi /2]`$ cannot be expanded by more than $`2\frac{1+\eta }{1\eta }2(1+ϵ)`$ under $`\mathrm{\Psi }`$, since that is the maximum of $`\frac{\mathrm{\Psi }^{}(\varphi _1)}{\mathrm{\Psi }^{}(\varphi _2)}`$ for $`\varphi _1,\varphi _2[\pi /2,\pi /2]`$. Now, since $`\nu ([\pi /2,\pi /2])=1/2`$, we have
$$\frac{\nu (A[\pi /2,\pi /2])}{\nu ([\pi /2,\pi /2])}2\nu (A),$$
which yields
$$\frac{\nu (\mathrm{\Psi }(A)U)}{\nu (U)}2(1+ϵ)2\nu (A)=4(1+ϵ)\nu (A).$$
Our second main lemma is as follows: let $`\tau >0`$. For any $`t>0`$, define
$$I(t):=\{0\theta <2\pi :d(i.g_t,i.g_tr_\theta )3\tau \}.$$
Fix $`\kappa >0`$. For any $`\theta I(t)`$, define
$$U_\theta :=\{0\varphi <2\pi :D_{t,\tau }(\varphi )>1\kappa ,(\mathrm{\Psi }_{t,\tau }(\varphi )+\theta )I(t+\tau )\}.$$
Let $`L_\theta =\nu (U_\theta )`$. The following lemma is proved in . For notational convenience, we write drop the subscripts for $`T`$ and $`S`$.
###### Lemma 4.6.
There is a constant $`c^{\prime \prime }>0`$ such that for all $`\kappa >0`$, there is a $`\tau >0`$ such that for all $`\theta I(t)`$, the map $`\mathrm{\Psi }|_{L_\theta }`$ is a diffeomorphism onto its image, and, making the subsitution $`\psi =\mathrm{\Psi }(\varphi )`$, we have
$$c^{\prime \prime }\nu (L_\theta )_{L_\theta }𝑑\varphi =_{\mathrm{\Psi }(L_\theta )}|\frac{d\varphi }{d\psi }|𝑑\psi .$$
Proof:, page 467, Lemma 7.6.
Remarks:
* $`d\theta `$ denotes $`d\nu (\theta )`$, and since we have normalized $`\nu `$ to be a probability measure, we do not need to divide by $`2\pi `$.
* In , the sets $`I(t)`$ are defined by $`I(t):=[\rho e^t,\rho e^t]`$ for some positive constant $`\rho `$, and they require $`\rho `$ to be large enough so that the diameter of the set $`J(t)=\{i.g_tr_\theta :\theta I(t)\}`$ is at least $`2\tau `$. By hyperbolic geometry, our sets $`I(t)`$ are of this form, and they obviously satisfy the required condition.
We have the following Corollary, also from :
###### Corollary 4.7.
Let $`f:SL(2,)`$ be a logsmooth $`SO(2)`$-invariant function. Fix $`\sigma >1`$. Let $`\kappa >0`$ be as in equation (2.8). Fix $`\tau `$ so that Lemma 4.6 holds. Then there is a $`c^{}>0`$, independent of $`\tau `$, such that
$$_{I(t+\tau )}f(g_{t+\tau }r_\theta )𝑑\theta c^{}\sigma _{I(t)}(A_\tau f)(g_tr_\theta )𝑑\theta .$$
Proof:, page 468, Lemma 7.7.
## 5. Large Deviations
In this section, we prove the key technical lemma for our main large deviations result Theorem 2.3. We assume some familiarity with the theory of conditional expectation. Excellent references include .
###### Proposition 5.1.
Let $`\{\tau _i\}_{i=0}^{\mathrm{}}`$ be a sequence of positive real-valued random variables on a probability space $`(\mathrm{\Omega },,P)`$. Let $`\{_i\}`$ be af filtration of $``$ such that for all $`i`$, $`\tau _i_i`$, i.e., $`\tau _i`$ is $`_i`$-measurable. Suppose there exist positive random variables $`\eta ,\xi `$ with $`E\eta <E\xi `$, $`E\xi >0`$, and a real number $`\theta _0>0`$ such that, for all $`0\theta <\theta _0`$:
1. $`E(e^{\theta \tau _{2i}}|_{2i1})E(e^{\theta \eta })`$
2. $`E(e^{\theta (\tau _{2i1}+\tau _{2i})}|_{2i1})E(e^{\theta \xi })`$
Let $`T_n=_{i=0}^n\tau _i`$. Let
$$X(t)=\{\begin{array}{cc}1\hfill & T_{2i1}t<T_{2i}\hfill \\ 0\hfill & \text{otherwise}\hfill \end{array}$$
Then, $`\lambda >E\eta /E\xi `$, there is a $`0<\gamma <1`$ such that for all $`T`$ sufficiently large
$$P\left(\frac{1}{T}_0^TX(t)𝑑t>\lambda \right)\gamma ^T.$$
For our application, $`\tau _{2i}`$ will be the time a random geodesic spends outside the compact set, and $`\tau _{2i1}`$ the time inside on the $`i`$th ‘sojourn’. $`\eta `$ is a stochastic upper bound for $`\tau _{2i}`$, and $`\xi `$ as a stochastic lower bound for the length of a ‘cycle’. Thus condition 1 should be thought of as a stochastic upper bound for the time spent outside and condition 2 a stochastic lower bound for the total time of a cycle. In our application, we will give a stronger, deterministic lower bound: in fact, we will construct our variables so that $`\tau _{2i1}>C`$, for some fixed $`C`$. This clearly implies condition 2, simply by taking $`\xi =C`$. For condition 1, we will use Theorem 2.1 to show that $`\tau _{2i}`$ cannot grow too large.
Proof:
Define $`f(\theta ):=E(e^{\theta \eta })`$ and $`g(\theta ):=E(e^{\theta \xi })`$.
Let $`N(T)=sup\nu \{k:T_{2k}T\}`$. Then,
(5.2) $`P\left({\displaystyle _0^T}X(t)𝑑t>\lambda T\right)`$ $``$ $`P\left({\displaystyle _0^T}X(t)𝑑t>\lambda T\text{ and }N(T)cT\right)+P\left(N(T)>cT\right)`$
$`=`$ $`P\left({\displaystyle \underset{i=0}{\overset{cT/2}{}}}\tau _{2i}>\lambda T\right)+P\left(N(T)>cT\right).`$
We estimate each of these terms in turn. Let $`n:=cT/2`$. Then, since
$$P\left(\underset{i=0}{\overset{cT/2}{}}\tau _{2i}>\lambda ^{}n\right)P\left(\underset{i=0}{\overset{cT/2}{}}\tau _{2i}>\lambda T\right)P\left(\underset{i=0}{\overset{cT/2}{}}\tau _{2i}>\lambda ^{}(n+1)\right)$$
to estimate the first term it suffices to estimate $`P\left(_{i=0}^n\tau _{2i}>\lambda ^{}n\right)`$, where $`\lambda ^{}=2\lambda /c`$. Now, for any $`\theta _0>\theta 0`$
(5.3) $`P\left({\displaystyle \underset{i=0}{\overset{n}{}}}\tau _{2i}>\lambda ^{}n\right)`$ $`=`$ $`P\left(e^{\theta _{i=0}^n\tau _{2i}}>e^{\theta \lambda ^{}n}\right)`$
$``$ $`e^{\theta \lambda ^{}n}E\left(e^{\theta _{i=0}^n\tau _{2i}}\right)`$
$``$ $`e^{\theta \lambda ^{}n}f(\theta )^n.`$
The last inequality follows from the first condition in our theorem, and the fact that each $`\tau _i`$ is $`_i`$-measurable. Since equation holds for any $`\theta _0>\theta 0`$, we have
$$P\left(\underset{i=0}{\overset{n}{}}\tau _{2i}>\lambda ^{}n\right)\underset{\theta _0>\theta 0}{inf}\left(f(\theta )e^{\theta \lambda ^{}}\right)^n.$$
Let $`\lambda ^{}>E\eta `$. Then, letting
$$F(\theta )=f(\theta )e^{\theta \lambda ^{}}=E\left(e^{\theta (\eta \lambda ^{})}\right),$$
we get $`F(0)=1`$,
$$F^{}(\theta )=E\left(e^{\theta (\eta \lambda ^{})}(\eta \lambda ^{})\right).$$
This implies that
$$F^{}(0^+)=E\left(\eta \lambda ^{}\right)<0.$$
Thus, there is a $`0\theta _1\theta _0`$ such that $`F(\theta _1):=\gamma ^{}<1`$. Plugging this into equation (5.3) yields the estimate for our first term.
To estimate the second term, let $`\xi _i=\tau _{2i1}+\tau _{2i}`$. Fix $`c>1/E\xi `$. By a similar argument to that above, we obtain
$$P\left(N(T)>cT\right)\underset{0\theta \theta _0}{inf}e^{\theta T}g(\theta )^{cT}.$$
Let
$$G(\theta )=g(\theta )e^{\theta /c}=E\left(e^{\theta (\xi 1/c)}\right).$$
Once again, as above, we obtain
$$G^{}(0^+)=E\left(\xi 1/c\right)0.$$
Thus, there exists $`\theta _2`$ with $`G(\theta _2)=\gamma ^{\prime \prime }<1`$, and so we have our desired estimate.∎
###### Corollary 5.4.
With notation as above,
$$\underset{T\mathrm{}}{lim\; sup}\frac{1}{T}_0^TX(t)𝑑t\lambda $$
with probability 1 for all $`\lambda >E\eta /E\xi `$.
In order to prove this corollary, we need the following technical lemma:
###### Lemma 5.5.
Let $`0<\gamma <1`$. Let $`U:^+^+`$ be such that for all sequences $`\{a_n\}_{n=0}^{\mathrm{}}`$ with $`_{n=0}^{\mathrm{}}\gamma ^{a_n}`$ convergent,
$$\underset{n\mathrm{}}{lim\; sup}U(a_n)c,$$
for some $`c>0`$. Then
$$\underset{T\mathrm{}}{lim\; sup}U(T)c.$$
Proof: We proceed by contradiction. Suppose $`lim\; sup_T\mathrm{}U(T)>c`$. Then, there is a sequence of time $`t_n`$, $`t_n\mathrm{}`$, such that $`U(t_n)>c`$. Take a subsequence $`t_{n_k}`$, where $`n_k`$ is such that $`t_n>k`$ for all $`nn_k`$. Such a subsequence exists since $`t_n`$ diverges. Now, letting $`a_k=t_{n_k}`$, note that $`a_k>k`$, so $`_{k=0}^{\mathrm{}}\gamma ^{a_k}`$ is convergent. So $`lim\; sup_k\mathrm{}U(a_k)c`$. But by definition, $`U(a_k)>c`$ for all $`k`$. This is a contradiction.∎
We now proceed with the proof of Corollary 5.4. Let $`U(T)=\frac{1}{T}_0^TX(t)𝑑t`$. Let $`\gamma `$ be as in the conclusion of Proposition 5.1. Then, for any sequence $`a_n`$ we have
$$\underset{n=0}{\overset{\mathrm{}}{}}P\left(U(a_n)>\lambda \right)\underset{n=0}{\overset{\mathrm{}}{}}\gamma ^{a_n}.$$
Thus, if $`_{n=0}^{\mathrm{}}\gamma ^{a_n}`$ converges, by the Borel-Cantelli lemma, $`lim\; sup_n\mathrm{}U(a_n)\lambda `$, with probability one. Applying Lemma 5.5, we have our result. ∎
## 6. Proofs of main results
We fix the following notation for the rest of this section: fix a $`\delta >0`$, and fix $`qQ`$. For $`hSL(2,)`$, we define $`V_{\delta ,q}(h)=V_\delta (hq)`$. $`V_{\delta ,q}`$ is $`K`$-invariant, and thus can be viewed as a function on $`^2`$. We define the sets $`C_{\delta ,l}(q)=\{z^2:V_q(z)l\}`$. In the rest of this section we work in this $`^2`$, identifying $`g_tr_\theta qQ`$ with $`i.g_tr_\theta ^2`$. For notational convenience, we drop the subscripts $`\delta `$ and $`q`$, and write $`V`$ and $`C_l`$ for $`V_{\delta ,q},C_{\delta ,l}(q)`$. Furthermore, all distances are measured in $`^2`$.
Proof of Theorem 2.1: Let $`qC_l`$. We want to estimate the measure of the sets
$$B^{}(T,l,q)=\{\theta :i.g_tr_\theta C_l,0tT\}.$$
For technical reasons, we will instead study the sets
$$B(T):=B(T,l,q)=\{\theta :\varphi B^{}(T,l,q)\text{ such that }d(i.g_Tr_\theta ,i.g_Tr_\varphi )3\tau \},$$
where we will specify $`\tau `$ shortly. Note that, by definition, and logsmoothness of $`V`$, there is some $`a_\tau 1`$ such that $`\varphi B(T,l,q)`$ implies that $`V(i.g_tr_\varphi )>l/a_\tau =:l^{\prime \prime }`$ for all $`0tT`$.
Let $`B_{n\tau }=B(n\tau ,l,q)`$, and $`p_{n\tau }=\nu (B_{n\tau })`$. We have
$$l^{\prime \prime }p_{n\tau }_{B_{n\tau }}V(i.g_{n\tau }r_\theta )d\theta =:D_{n\tau }.$$
Our main lemma is as follows
###### Lemma 6.1.
For all $`\delta >0`$, and all $`\tau `$ sufficiently large there are constants $`c=c(\delta ),b=b(\tau ,\delta )`$, so that
$$D_{n\tau }ce^{(1\delta )\tau }+b$$
Proof: Note that by definition $`B_{n\tau }`$ is a union of arcs of the form $`I(n\tau )`$, and as such, both Lemma 4.6 and Corollary 4.7 apply (we are also using the fact that $`V`$ is logsmooth). Fix $`\sigma >1`$, and let $`\tau `$ be such that we can apply corollary 4.7. Setting $`c=\sigma c^{}\stackrel{~}{c}`$ and $`b=\sigma c^{}\stackrel{~}{b}`$, we obtain our result.∎
Let $`\tau `$ be large enough so that we can apply Lemma 6.1. Proceeding as in the proof of Proposition 3.1 part (1), we obtain
$$p_{n\tau }(q)\frac{V_\delta (q)}{l^{\prime \prime }}\left(ce^{(1\delta )\tau }+\frac{b}{l^{\prime \prime }}\right)^n.$$
Let $`\tau _0>0`$ be such that $`ce^{(1\delta )\tau _0}<1`$.
Let $`t>\tau _0`$, and let
$$l_0a\underset{\tau _0\tau 2\tau _0}{sup}\frac{b}{(1ce^{1\delta )\tau })}.$$
Let
$$a=\underset{\tau _0\tau 2\tau _0}{sup}a(\tau ),$$
and set $`l^{}=l/a`$.
Let
$$\delta ^{}=\delta +\underset{\tau _0\tau 2\tau _0}{sup}\frac{1}{\tau }\mathrm{ln}\left(c+\frac{b}{l^{}}e^{(1\delta )\tau }\right).$$
It is easy to check that $`\delta ^{}<1`$, and that it is decreasing as a function of $`l`$. There is some $`\tau _0\tau 2\tau _0`$ and $`n`$ such that $`t=n\tau `$. We have
$$p_{n\tau }\frac{V_\delta (q)}{l^{}}\left(ce^{(1\delta )\tau }+\frac{b}{l^{}}\right)^n.$$
Rewriting this, we obtain
$$p_ta\frac{V_\delta (q)}{l}e^{(1\delta ^{})t}.$$
Proof of Theorem 2.2: Let $`qQ`$. Consider the circle of radius $`S+T`$, $`\{i.g_{S+T}r_\theta :0\theta <2\pi \}`$. Given $`T>0`$, we want to show that the set $`B=B_{S,T}(q,l)=\{\theta :i.g_tr_\theta C_l,StS+T\}`$ has exponentially small measure in $`T`$ for sufficiently large $`l,S,`$ and $`T`$.
Given $`\theta _0B`$, let $`z_0=i.g_Sr_{\theta _0}`$, and consider the circle $`\{i.g_Tr_\varphi g_Sr_{\theta _0}:0\varphi <2\pi \}`$ of radius $`T`$ around it. By Theorem 2.1, we know that for most (the complement is exponentially small in $`T`$) directions $`\varphi `$ on this circle, $`i.g_tr_\varphi g_Sr_{\theta _0}C_l`$ for some $`t<T`$.
Our idea is as follows: there is a small neighborhood $`U`$ of $`\theta _0`$ such that each geodesic trajectory $`\{i.g_tr_\theta \}_{t=0}^{S+T},\theta U`$ is closely shadowed by a piecewise geodesic of the form $`\gamma =\gamma _{\theta _0,\varphi }`$, $`\theta =\mathrm{\Psi }_{S,T}(\varphi )`$, where
$$\gamma (t)=\{\begin{array}{cc}i.g_tr_{\theta _0}\hfill & 0tS\hfill \\ i.g_{tS}r_\varphi g_Sr_{\theta _0}\hfill & S<t<S+T\hfill \end{array}$$
By closely shadowed, we mean that $`d(\gamma (t),i.g_tr_\theta )`$ is small for all $`t`$. Thus, if $`V(\gamma (t))l`$ for some $`SS+T`$, we have that $`V(i.g_tr_\theta )\stackrel{~}{l}`$, for some $`\stackrel{~}{l}>l`$.
Now, for all but a small set of $`\varphi `$, $`\gamma (t)C_l`$ for some $`StS+T`$. Thus, in a small neighborhood of $`\theta _0`$, we have a (large-proportioned) collection of angles which are not in $`\stackrel{~}{B}=B_{S,T}(q,\stackrel{~}{l})`$.
To make this rigorous, fix $`ϵ,\delta >0`$, and let $`A:=A(\theta _0)=B_{0,T}(z_0,l)=\{\varphi :V(i.g_tr_\varphi g_Sr_{\theta _0})>l,0tT\}`$. Let $`d`$ be the maximum thickness of a hyperbolic triangle. By logsmoothness of $`V`$, there is a $`b1`$ such that $`d(z_1,z_2)d`$ implies $`\frac{V(z_1)}{V(z_2)}b`$. Let $`l_1=bl_0`$, where $`l_0`$ is as in the conclusion of Theorem 2.1. For any $`l`$, let $`\widehat{l}=l/b`$.
Setting $`\delta ^{\prime \prime }=\delta ^{}(\widehat{l},\delta )`$ and $`l^{}=\widehat{l}/a=l/ba`$, where $`a=a(\delta )`$ is as in Theorem 2.1 yields
$$\nu (A)\frac{V(z_0)}{l^{}}e^{(1\delta ^{\prime \prime })t}.$$
Let $`T,S`$ be large enough so that we can apply Lemma 4.3 with $`ϵ`$, and $`t_2=T`$, $`t_1=S`$. This gives a neighborhood $`U`$ of $`\theta _0`$ in which the proportion of angles
$$\frac{\nu (\mathrm{\Psi }(A)U)}{\nu (U)}4(1+ϵ)\underset{\theta [0,2\pi )}{sup}\frac{V(i.g_Sr_\theta )}{l^{}}e^{(1\delta ^{\prime \prime })t}.$$
Now, by the thinness of triangles in hyperbolic geometry, given $`\theta U`$, there is a $`\varphi `$ with $`\mathrm{\Psi }(\varphi )=\theta `$ such that
$$d(i.g_tr_\theta ,\gamma _{\theta _0,\varphi }(t))d$$
for $`t[0,S+T]`$.
For all $`\varphi A`$, we have
$$\gamma _{\theta _0,\varphi }(t_0)C_{\widehat{l}},$$
for some $`S+T>t_0>S`$. Thus,
$$g_{t_0}r_\theta qC_l,$$
so $`\theta B`$.
Thus, given any $`\theta _0B`$, we have produced a neighborhood $`U`$ s.t.
$$\frac{\nu (B^CU)}{\nu (U)}>14(1+ϵ)\underset{\theta [0,2\pi )}{sup}\frac{V(i.g_Sr_\theta )}{l^{}}e^{(1\delta ^{\prime \prime })t}.$$
To complete the proof, we need the following standard lemma (see, for example, ):
###### Lemma 6.2.
Let $`B[0,2\pi )`$ be a measurable set such that for all $`bB`$, there is a $`\delta _b>0`$, so that $`U_b=[b\delta _b,b+\delta _b][0,2\pi )`$ satisfies
$$\frac{\nu (U_bB)}{\nu (U_b)}<ϵ.$$
Then
$$\nu (B)2ϵ.$$
Applying the lemma to our set $`B`$, we obtain
$$\nu (B)8(1+ϵ)\underset{\theta [0,2\pi )}{sup}\frac{V(i.g_Sr_\theta )}{l^{}}e^{(1\delta ^{\prime \prime })t}.$$
Proof of Theorem 2.3: Our strategy is as follows: Given a direction $`\theta `$, consider the succesive departures and returns of the geodesic trajectory $`\{i.g_tr_\theta \}_{t0}`$ to the compact set $`C_l`$. Theorem 2.1 implies that the probability any departure is long is small, and thus, we can try and apply Proposition 5.1, to the ‘random variables’ given by the length of sojourns inside and outside the compact set.
We proceed as follows: Let $`d`$ be as in the proof of Theorem 2.2. Fix $`\delta >0`$, and let $`l_0`$ be as in Theorem 2.1. Let $`l>l_0`$ be such that $`d(C_l^c,C_{l_0})>2d`$, and define $`C=C(l)=d(C_l^c,C_{l_0})2d`$. Define $`t_0(\theta )=0`$, and set
$$t_{2n}(\theta )=inf\{t>t_{2n1}:\varphi \text{ such that }d(i.g_tr_\varphi ,i.g_tr_\theta )<d,i.g_tr_\varphi C_{l_0}\}$$
and
$$t_{2n+1}(\theta )=inf\{t>t_{2n}:\varphi \text{ such that }d(i.g_tr_\varphi ,i.g_tr_\theta )<d,i.g_tr_\varphi C_l\},$$
for $`n0`$. Define $`\tau _i(\theta )=t_it_{i1}`$. Now fix $`C^{}>0`$, and define auxiliary functions $`\tau _i^{}`$ by
$$\tau _{2i}^{}:=\{\begin{array}{cc}0\hfill & \tau _{2i}C^{}\hfill \\ \tau _{2i}\hfill & \tau _{2i}>C^{}\hfill \end{array}$$
and
$$\tau _{2i1}^{}:=\{\begin{array}{cc}\tau _{2i1}+\tau _{2i}\hfill & \tau _{2i}C^{}\hfill \\ \tau _{2i1}\hfill & \tau _{2i}>C^{}\hfill \end{array}$$
Define $`t_i^{}:=_{j=1}^i\tau _j^{}`$. Note that $`\tau _{2i}+\tau _{2i1}=\tau _{2i}^{}+\tau _{2i1}^{}`$, so $`t_{2i}^{}=t_{2i}`$.
$$X(t)=\{\begin{array}{cc}1\hfill & t_{2i1}^{}t<t_{2i}^{}\hfill \\ 0\hfill & \text{otherwise}\hfill \end{array}$$
The $`t_i`$’s should be thought of as the entry and departure times of the trajectory $`\{i.g_tr_\theta \}_{t0}`$ from our compact sets. For technical reasons, they are defined as the first time when a *nearby* trajectory leaves a larger compact set ($`C_l`$), or re-enters a smaller one ($`C_{l_0}`$). The $`\tau _{2i}`$’s measure time spent after departing the larger set before returning to the smaller, and the $`\tau _{2i1}`$’s measure the time spent after returning to the smaller before departing the larger. The auxiliary $`\tau _i^{}`$’s are defined to exclude short ($`C^{}`$) sojourns. Thus, if $`t_{2i}^{}<t<t_{2i+1}^{}`$ we are within distance $`C^{}`$ of the larger compact set, and thus still within a compact set. Thus, for all such $`t`$, $`V(i.g_tr_\theta )\stackrel{~}{l}=\stackrel{~}{l}(C^{})`$, i.e.,
$$X(t)\chi _{C_{\stackrel{~}{l}}^c}(i.g_tr_\theta ).$$
Thus, it suffices to show that for all $`\lambda >0`$ there is a $`C^{}>0`$ so that $`\nu \left(\frac{1}{T}_0^TX(t)𝑑t>\lambda \right)`$ decays exponentially in $`T`$.
We will apply Proposition 5.1 to $`\{\tau _i^{}\}_{i0}`$, with $`\mathrm{\Omega }=S^1`$, $``$ the standard $`\sigma `$-algebra, and $`P=\nu `$ the Haar measure. Let $`_n=\sigma (t_1^{},\mathrm{},t_n^{})`$ be the $`\sigma `$-algebra generated by $`t_1^{},\mathrm{},t_n^{}`$. Clearly $`\tau _n^{}`$ is $`_n`$-measurable. Note that since $`d(C_l^c,C_{l_0})=2d+C`$,
$$\tau _{2i+1}^{}\tau _{2i+1}=t_{2i+1}t_{2i}>C$$
(since $`i.g_{t_{2i}}r_\theta `$ is within distance $`d`$ of $`C_{l_0}`$ and $`i.g_{t_{2i+1}}r_\theta `$ is within distance $`d`$ of $`C_l^c`$). Thus, condition 2 of the proposition is satisfied, with $`\xi =C`$
It remains to check condition 1. We will show
(6.3)
$$\nu (\tau _{2i}^{}>t|_{2i1})a_1e^{a_2t}$$
for some $`a_1,a_2>0`$ and all $`t>C^{}`$.
By definition, $`t_i^1(x)`$ is an interval, and thus $`_i`$ is generated by these intervals. Fix $`\theta `$. Let
$$I_n(\theta )=\{\varphi :t_i^{}(\varphi )=t_i^{}(\theta )\text{ for all }0in\}.$$
It suffices to show that there are $`a_1,a_2>0`$ such that, for all $`t>C^{}`$,
(6.4)
$$\nu \left(\tau _{2n}^{}>t|I_{2n1}(\theta )\right)a_1e^{a_2t}.$$
Fix $`\varphi I_{2n1}(\theta )`$ with $`\tau _{2i}(\varphi )>t`$. Let $`z_\varphi =i.g_{t_{2i1}}r_\varphi `$. Note that this point is within distance $`d`$ of the boundary of $`C_l^c`$, thus it is both outside $`C_{l_0}`$ and still contained within a compact set. Consider the circle of radius $`t`$ around $`z_\varphi `$ and the associated map $`\mathrm{\Psi }=\mathrm{\Psi }_{t,t_{2i1}}`$ back to the circle at $`i`$.
Let $`A=A(\varphi )=\{\theta :i.g_sr_\theta g_{t_{2i1}}r_\varphi C_{l_0},0st\}`$. Since $`z_\varphi `$ is contained in a compact set, we can pick $`c_1=c_1(l),c_2=c_2(l_0)`$ independent of $`\varphi `$ such that
$$\nu (A)c_1e^{c_2t},$$
for all $`t`$ sufficiently large. By picking $`C^{}`$ large enough, we can get this to hold for all $`t>C^{}`$. For the rest of this section, let $`t>C^{}`$.
Applying lemma 4.3 with $`ϵ=1`$, we obtain a neighborhood $`U`$ of $`\varphi `$, with
$$\frac{\nu (\mathrm{\Psi }AU)}{\nu (U)}16c_1e^{c_2t}.$$
Note that $`U\chi _{2n1}(\theta )`$, since $`0tt_{2i1},\theta ^{}U`$,
$$d(i.g_tr_\theta ^{},i.g_tr_\theta )<d.$$
Finally, observe that $`\tau _{2i}(\theta ^{})<t`$ for all $`\theta ^{}U\mathrm{\Psi }(A(\varphi ))`$, since $`\theta ^{}\mathrm{\Psi }(A)`$ implies that there is some $`0st`$ such that $`V(i.g_sr_\psi g_{t_{2i1}}r_\varphi )l_0`$, with $`\mathrm{\Psi }(\psi )=\theta ^{}`$. Thus, since $`d(i.g_sr_\psi g_{t_{2i1}}r_\varphi ,i.g_{s+t_{2i1}}r_\theta ^{})d`$, we have $`V(i.g_{s+t_{2i1}}r_\theta ^{})l`$, and thus $`\tau _{2i}(\theta ^{})<t`$. Once again applying Lemma 6.2, we obtain equation (6.4), with $`a_1=32c_1`$ and $`a_2=c_2`$.
Let $`C^{}`$ be large enough so that $`a_1e^{a_2C^{}}<1`$. Let $`\eta `$ be a non-negative function on $`S^1`$ such that
$$\nu \{\theta :\eta (\theta )=0\}=a_1e^{a_2C^{}}$$
and
$$\nu \{\theta :\eta (\theta )>t\}=a_1e^{a_2t}$$
for all $`t>C^{}`$. Condition 1 is then clearly satisfied, since $`\nu (\tau _{2i}^{}=0)a_1e^{a_2C^{}}`$, and $`\nu (\tau _{2i}^{}>t)a_1e^{a_2t}`$ for $`t>C^{}`$. Now, note that $`E(\eta )=\frac{a_1}{a_2}e^{a_2C^{}}`$, and $`E(\xi )=C`$. Thus, by enlarging $`C^{}`$, we can make $`E(\eta )/E(\xi )`$ arbitrarily small. Precisely for any $`\lambda >0`$, take $`C^{}`$ so that
$$E(\eta )/E(\xi )\lambda .$$
Then, setting $`\stackrel{~}{l}=\stackrel{~}{l}(C^{})`$ and applying Proposition 5.1, we obtain that there is a $`\gamma <1`$ so that
$$\nu \{\theta :\frac{1}{T}|\{0tT:g_tr_\theta qC_{\stackrel{~}{l}}\}|>\lambda \}\nu \left(\frac{1}{T}_0^TX(t)𝑑t>\lambda \right)\gamma ^T,$$
for all $`T`$ sufficiently large. ∎
Proof of Corollary 2.4: Applying Corollary 5.4 to $`X(t)`$, we obtain our result.∎
Acknowledgements: I would like to thank my advisor, Professor Alex Eskin, for his guidance throughout this project. I would also like to thank Professors Howard Masur, Giovanni Forni, Steven Lalley and Krishna Athreya for valuable discussions, and my colleague Matthew Day for help with hyperbolic trigonometry. Thanks are also due to the anonymous referee, whose detailed comments and suggestions greatly improved this paper.
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# Clustering by mixing flows
## Abstract
We calculate the Lyapunov exponents for particles suspended in a random three-dimensional flow, concentrating on the limit where the viscous damping rate is small compared to the inverse correlation time. In this limit Lyapunov exponents are obtained as a power series in $`ϵ`$, a dimensionless measure of the particle inertia. Although the perturbation generates an asymptotic series, we obtain accurate results from a Padé-Borel summation. Our results prove that particles suspended in an incompressible random mixing flow can show pronounced clustering when the Stokes number is large and we characterise two distinct clustering effects which occur in that limit.
This letter describes the dynamics of particles suspended in a randomly moving incompressible fluid which we assume to be mixing: any given particle uniformly samples configuration space. At first sight, it seems as if the particles suspended in an incompressible mixing flow should become evenly distributed. This indeed happens if the particles are simply advected by the fluid. However, it has been noted Max87 that when the finite inertia of the suspended particles is significant, the particles can show a tendency to cluster.
The current understanding of this remarkable phenomenon refers to a dimensionless parameter termed the Stokes number, $`\mathrm{St}=1/(\gamma \tau )`$, where $`\gamma `$ is the rate at which the particle velocity is damped relative to that of the fluid due to viscous drag, and $`\tau `$ is the correlation time of the velocity of the fluid. There is a consensus Fes94 ; Hog01 ; Sig02 ; Bec03 ; Gam04 that clustering is most pronounced when $`\mathrm{St}`$ is of order unity.
In this letter we argue that strong clustering can occur when St is large. We show that different clustering mechanisms compete at large values of St and quantify under which circumstances clustering occurs. Before describing our results and outlining how they are derived, we briefly summarise previous theoretical work on the clustering of inertial particles in turbulent flows.
This effect was first discussed by Maxey Max87 : he approximated the inertial particle dynamics by advection in a ‘synthetic’ velocity field which was obtained as a perturbation of the velocity field of the fluid, $`𝒖(𝒓,t)`$. Maxey showed that this synthetic velocity field has negative divergence when the vorticity of $`𝒖(𝒓,t)`$ is high or its strain-rate low, and predicted that particles would have low concentrations in regions of high vorticity due to this ‘centrifuge effect’. This effect has been demonstrated in direct numerical simulation of particles suspended in a fully-developed turbulent flow Wan93 ; Hog01 . The theoretical work of Maxey and experimental work on turbulent flows Fes94 has emphasised instantaneous correlations between vortices and particle-density fluctuations.
Later work has adapted results on the density statistics and Lyapunov exponents of purely advective flows obtained in LeJ85 ; Ber00 : Elperin Elp96 suggested combining these results with Maxey’s synthetic velocity field to obtain results for inertial particles; a similar approach was used in Pin99 ; Bal01 ; Fal02 . These results are not applicable at large St, because the perturbation of the velocity field need not be small when inertial effects are important.
An alternative viewpoint arises from work of Sommerer and Ott Som93 , who describe patterns formed by particles floating on a randomly moving fluid. They characterise the patterns in terms of their fractal dimension and suggest that the fractal dimension can be obtained from ratios of Lyapunov exponents of the particle trajectories using a formula proposed by Kaplan and Yorke Kap79 .
The argument in Som93 extends to particles suspended in turbulent three-dimensional incompressible flows. Consider the Lyapunov exponents $`\lambda _1>\lambda _2>\lambda _3`$. They are rate constants defined in terms of the time dependence of, respectively, the length $`\delta r`$ of a small separation between two trajectories, the area $`\delta 𝒜`$ of a parallelogram spanned by two separation vectors and the volume $`\delta 𝒱`$ of a parallelepiped spanned by a triad of separations:
$`\lambda _1`$ $`=`$ $`\underset{t\mathrm{}}{lim}t^1\mathrm{log}_\mathrm{e}(\delta r)`$
$`\lambda _1+\lambda _2`$ $`=`$ $`\underset{t\mathrm{}}{lim}t^1\mathrm{log}_\mathrm{e}(\delta 𝒜)`$
$`\lambda _1+\lambda _2+\lambda _3`$ $`=`$ $`\underset{t\mathrm{}}{lim}t^1\mathrm{log}_\mathrm{e}(\delta 𝒱).`$ (1)
The Kaplan-Yorke estimate for the fractal dimension in a three-dimensional incompressible flow is determined by the dimensionless quantity (‘dimension deficit’)
$$\mathrm{\Delta }=(\lambda _1+\lambda _2+\lambda _3)/|\lambda _3|.$$
(2)
When $`\mathrm{\Delta }>0`$, the Kaplan-Yorke estimate of the dimension is $`d_\mathrm{H}=3\mathrm{\Delta }`$, and $`d_\mathrm{H}=3`$ if $`\mathrm{\Delta }0`$. Clustering effects are significant if the fractal dimension is significantly lower than the dimension of space. This proposition provides a strong motivation to study the Lyapunov exponents of the problem.
A third mechanism for clustering is the following: nothing prevents the infinitesimal volume element $`\delta 𝒱`$ from collapsing to zero for an instant of time. These events correspond to ‘caustics’, where faster moving particles overtake slower ones. Caustics are associated with the density of particles on a surface becoming very high, facilitating the aggregation of suspended particles. This mechanism was recently proposed as a cause of clustering of inertial particles Wil05 , and is also mentioned briefly in Fal02 . The significance of this effect is determined by the rate $`J`$ at which the infinitesimal volume element goes through zero for a given triplet of nearby trajectories.
Which of these three mechanisms is most important? Maxey’s centrifuge effect is weak at small $`\mathrm{St}`$, where the particles are simply advected. There is a consensus that the effect is also weak for large St, because the vortices do not persist for a sufficiently long time to be effective, implying that significant clustering is only observed when $`\mathrm{St}1`$. However, there is at present no understanding of what happens at large values of St. In the following we describe quantitative results for the Lyapunov exponents $`\lambda _j`$, for the dimension deficit $`\mathrm{\Delta }`$ and for the rate of caustic formation $`J`$: these are summarised in Fig. 1 a \- c.
Our results show that in order to understand the clustering effect it is necessary to consider not only the Stokes number, but an additional dimensionless parameter, $`\kappa `$, defined below. We infer that strong clustering can occur at large Stokes numbers. Two distinct mechanisms compete (clustering onto fractal sets versus clustering onto caustics in an otherwise homogeneous background) and dominate in different regions of the parameter space.
We model the particles suspended in the fluid flow by the equation of motion
$$\ddot{𝒓}=\gamma \left(𝒖(𝒓,t)\dot{𝒓}\right)$$
(3)
where $`𝒓=(r_1,r_2,r_3)`$ denotes the position of a particle. Eq. (3) is appropriate for non-interacting spherical particles when the Reynolds number of the flow referred to the particle diameter is small. It is assumed that the radius of the particle and the molecular mean free path of the fluid are sufficiently small. Stokes’s formula gives the damping rate $`\gamma =6\pi a\rho _\mathrm{f}\nu /m`$ where $`\nu `$, $`\rho _\mathrm{f}`$ are respectively the kinematic viscosity and density of the fluid, and $`a`$, $`m`$ are the radius and mass of the particle. Effects due to the inertia of the displaced fluid are neglected. This is justified when the density of the suspended particles is large compared to that of the fluid. We also assume that Brownian diffusion of the particles is negligible.
We now discuss the dimensionless parameters of the problem: the velocity field is assumed to be characterised by its typical velocity $`u=\sqrt{𝒖^2}`$, by a correlation length $`\xi `$ and a correlation time $`\tau `$. In addition, the interaction of the fluid with the particles is determined by the damping rate $`\gamma `$. From these four quantities we can form two independent dimensionless groups: a dimensionless velocity, $`\kappa =u\tau /\xi `$, and the dimensionless damping $`\omega =\gamma \tau `$ (so that $`\text{St}=\omega ^1`$). The parameter $`\kappa `$ has been termed ‘Kubo number’ Bri74 . It has not been considered before in this context. We argue that it cannot be large if $`𝒖(𝒓,t)`$ is to be a satisfactory model for a solution of the Navier-Stokes equations: $`\tau \xi /u`$ since disturbances in the fluid velocity field $`𝒖(𝒓,t)`$ are transported by $`𝒖(𝒓,t)`$ itself.
Consider now the particular case of fully-developed turbulence. In this case, the velocity field exhibits a power-law energy spectrum, with upper and lower cutoffs Fri97 . The smaller length scale is the Kolmogorov length, which is the size of the smallest vortices generated by the turbulence. It is given by $`(\nu ^3/\epsilon )^{1/4}`$, where $`\epsilon `$ is the rate of dissipation per unit mass of fluid. The Kolmogorov length corresponds to the correlation length $`\xi `$ in our theory. The corresponding typical velocity $`u`$ and correlation time $`\tau `$ are also determined solely by the same two parameters, $`\epsilon `$ and $`\nu `$, implying that $`\kappa 1`$ for fully developed turbulence. In other situations $`\kappa `$ can be small.
We now turn to a summary of our results and outline how they were derived (details will be published elsewhere). Linearising the equations of motion (3) gives
$`\delta \dot{𝒑}`$ $`=`$ $`\gamma \delta 𝒑+𝐅(t)\delta 𝒓,\delta \dot{𝒓}=\delta 𝒑/m`$ (4)
where $`𝒑=m\dot{𝒓}`$ is the particle momentum and $`𝐅(t)`$ is matrix of force gradients:
$$F_{\mu \nu }(t)=\gamma m\frac{u_\mu }{r_\nu }(𝒓(t),t).$$
(5)
We take three trajectories displaced relative to a reference trajectory by small increments $`(\delta 𝒓_\mu ,\delta 𝒑_\mu )`$, with $`\mu =1,2,3`$. We introduce a triplet of orthogonal unit vectors $`𝐧_\nu (t)`$ such that $`𝐧_1(t)`$ is oriented along $`\delta 𝒓_1(t)`$, and $`𝐧_2(t)`$ lies in the plane spanned by $`(\delta 𝒓_1(t),\delta 𝒓_2(t))`$. This determines $`𝐧_3(t)`$ up to a sign which is fixed by requiring continuity. We write $`𝐧_\nu (t)=𝐎(t)𝐧_\nu (0)`$ and $`\delta 𝒑_\mu (t)=𝐑(t)\delta 𝒓_\mu (t)`$ where $`𝐎`$ is an orthogonal and $`𝐑`$ a general $`3\times 3`$ matrix. We define the elements of $`𝐅`$ and $`𝐑`$ transformed to the moving basis:
$$F_{\mu \nu }^{}(t)=𝐧_\mu (t)𝐅(t)𝐧_\nu (t),R_{\mu \nu }^{}(t)=𝐧_\mu (t)𝐑(t)𝐧_\nu (t)$$
(6)
and find the following equation of motion for $`𝐑^{}`$
$$\dot{𝐑^{}}=\gamma 𝐑^{}\frac{1}{m}𝐑_{}^{}{}_{}{}^{2}+[𝐑^{},𝐎^+\dot{𝐎}]+𝐅^{}.$$
(7)
The elements of $`𝐎^+\dot{𝐎}`$ are given by
$$𝐎^+\dot{𝐎}=\frac{1}{m}\left(\begin{array}{ccc}0\hfill & R_{21}^{}& \hfill R_{31}^{}\\ R_{21}^{}\hfill & 0& \hfill R_{32}^{}\\ R_{31}^{}\hfill & R_{32}^{}& \hfill 0\end{array}\right).$$
(8)
We find that the Lyapunov exponents are equal to the long-time average of the diagonal elements of $`𝐑^{}`$
$`\lambda _1`$ $`=`$ $`R_{11}^{}/m,\lambda _2=R_{22}^{}/m,\lambda _3=R_{33}^{}/m.`$ (9)
Eqs. (7) and (8) for $`𝐑^{}`$ can be simplified when the correlation time of the velocity field is sufficiently short, $`\omega 1`$, assuming that the amplitude of the random force is sufficiently small, $`\kappa 1`$. In this limit $`𝐅^{}`$ behaves as a white-noise signal, and (7) reduces to a system of Langevin equations. We label the dynamical variables by a single index $`i=3(\mu 1)+\nu `$ and scale the Langevin equations for $`R_i^{}`$ to dimensionless form
$$\mathrm{d}x_i=\left(x_i+ϵ\underset{j=1}{\overset{9}{}}\underset{k=1}{\overset{9}{}}V_{jk}^ix_jx_k\right)\mathrm{d}t^{}+\mathrm{d}w_i$$
(10)
Here $`t^{}=\gamma t`$, $`x_i=\sqrt{\gamma /D_1}R_i^{}`$, and $`\mathrm{d}w_i\mathrm{d}w_j=2D_{ij}\mathrm{d}t^{}`$. The elements $`D_{ij}`$ of the diffusion matrix $`𝐃`$ are given by
$$D_{ij}=\frac{1}{2}_{\mathrm{}}^{\mathrm{}}dtF_i^{}(t)F_j^{}(0).$$
(11)
The coefficients $`V_{jk}^i`$ are determined by the 2nd and 3rd terms on the rhs of (7). The dimensionless parameter
$$ϵ=D_{11}^{1/2}/(m\gamma ^{3/2})\kappa \omega ^{1/2}$$
(12)
is a measure of the inertia of the particles: it is proportional to $`a`$ and therefore to $`m^{1/3}`$. Thus we obtain all three Lyapunov exponents from the expectation values of variables in a system of Langevin equations. Earlier work has obtained the largest Lyapunov exponents for various problems using Langevin equations Hal65 ; Pit02 ; Meh04 .
The elements of $`𝐃`$ are determined by the fluctuations of the velocity field. We assume that the latter is incompressible, but for reasons explained below we add a small compressible component: $`𝒖=\mathbf{}𝑨+\mathbf{}\delta A_0`$. The fields $`A_\mu (𝒓,t),\mu =1,2,3`$ are taken to be homogeneous in space and time, and isotropic in space. Their correlations are determined by $`A_\mu (𝒓+𝑹,t_0+t)A_\nu (𝒓_0,t_0)=\delta _{\mu \nu }C(|𝒓𝒓^{}|/\xi ,|tt^{}|/\tau )`$. The field $`\delta A_0`$ is statistically independent of $`A_\mu `$, has the same correlation function, and in the end the limit $`\delta A_00`$ is taken.
The Langevin equations (10) are equivalent to a Fokker-Planck equation whose stationary solution $`P(𝒙)`$ determines the Lyapunov exponents. In the limit of $`ϵ0`$ the latter is Gaussian
$$P_0(𝒙)\mathrm{exp}(\frac{1}{2}𝒙𝐃^1𝒙)\mathrm{exp}[\mathrm{\Phi }_0(𝒙)].$$
(13)
This suggests transforming the Fokker-Planck operator so that its $`ϵ=0`$ limit is transformed into a harmonic oscillator. This is achieved by introducing $`Q(𝒙)=\mathrm{exp}[\mathrm{\Phi }_0(𝒙)/2]P(𝒙)`$. The steady-state Fokker-Planck equation can be written as $`(\widehat{H}_0+ϵ\widehat{H}_1)|Q)=0`$, where we have represented the function $`Q(𝒙)`$ by a ‘ket vector’ $`|Q)`$. The operator $`\widehat{H}_0`$ is the Hamiltonian for nine uncoupled harmonic oscillators
$$\widehat{H}_0=\underset{i=1}{\overset{9}{}}\widehat{a}_i^+\widehat{a}_i$$
(14)
where the $`\widehat{a}_i^+`$ and $`\widehat{a}_i`$ are, respectively, the creation and annihilation operators for the degree of freedom labelled by $`i`$ (satisfying $`[\widehat{a}_i,\widehat{a}_j^+]=\delta _{ij}\widehat{I}`$). The non-Hermitean perturbation $`\widehat{H}_1`$ can be expressed in terms of the eigenvalues $`\omega _i`$ of $`𝐃`$ and the elements $`J_{ij}`$ of an orthogonal matrix $`𝐉`$ satisfying $`𝐃=𝐉𝛀𝐉^1`$, with $`𝛀=\mathrm{diag}(\omega _i)`$:
$`\widehat{H}_1`$ $`=`$ $`{\displaystyle \underset{ijk}{}}H_{ijk}^{(1)}\widehat{a}_i^+(\widehat{a}_j^++\widehat{a}_j)(\widehat{a}_k^++\widehat{a}_k)`$
$`H_{ijk}^{(1)}`$ $`=`$ $`\sqrt{\omega _j\omega _k/\omega _i}{\displaystyle \underset{lmn}{}}V_{mn}^lJ_{il}J_{mj}J_{nk}.`$ (15)
Regularisation is needed since one eigenvalue vanishes in the limit of $`\delta A_00`$. We determine $`|Q)`$ by perturbation theory in $`ϵ`$. Given $`|Q)`$, the Lyapunov exponents are obtained as $`\lambda _1=\gamma ϵx_1`$, $`\lambda _2=\gamma ϵx_5`$, $`\lambda _3=\gamma ϵx_9`$, and
$`x_i`$ $`=`$ $`{\displaystyle \frac{1}{(\mathrm{\Phi }_\mathrm{𝟎}|Q)}}{\displaystyle \underset{j}{}}J_{ij}\sqrt{\omega _j}(\mathrm{\Phi }_\mathrm{𝟎}|\widehat{a}_j+\widehat{a}_j^+|Q)`$ (16)
where $`|\mathrm{\Phi }_0)`$ denotes the null eigenvector of $`\widehat{H}_0`$. From (16) we obtain series expansions in the form
$`\lambda _1/\gamma `$ $`=`$ $`3ϵ^229ϵ^4+564ϵ^6`$
$`14977ϵ^8+488784ϵ^{10}18670570ϵ^{12}+\mathrm{}`$
$`\lambda _2/\gamma `$ $`=`$ $`8ϵ^4459/2ϵ^6+14281/2ϵ^8`$
$`757273/3ϵ^{10}+361653709/36ϵ^{12}+\mathrm{}`$
$`\lambda _3/\gamma `$ $`=`$ $`3ϵ^29ϵ^4789/2ϵ^65787/2ϵ^8`$
$`895169/3ϵ^{10}101637719/36ϵ^{12}+\mathrm{}.`$
Note that only even powers of $`ϵ`$ contribute, and that all coefficients are rational numbers. Eq. (Clustering by mixing flows) is the main result of this letter. The expansion is valid in the underdamped limit $`\omega 1`$ when $`\kappa 1`$.
The coefficients in (Clustering by mixing flows) exhibit rapid growth typical of an asymptotic series Boy99 . We have attempted to sum the series (Clustering by mixing flows) using Padé-Borel summation Boy99
$$\lambda _j/\gamma \text{Re}_Cdt\mathrm{e}^t\underset{l=1}{\overset{\mathrm{l}_{\mathrm{max}}}{}}\frac{c_l^{(j)}}{l!}ϵ^{2l}$$
(18)
where $`c_l^{(j)}`$ are the coefficients of (Clustering by mixing flows) and $`l_{\mathrm{max}}=7`$ is the number of nonzero coefficients available for each $`\lambda _j`$. The sum in the integrand is approximated by Padé approximants Ben78 of order $`n`$, namely $`P_n^n`$ or $`P_{n+1}^n`$ with $`n[l_{\mathrm{max}}/2]`$. The integration path in (18) is taken to be a ray in the upper right quadrant in the complex plane.
Results of Padé-Borel summations of the series for $`\lambda _j`$ are shown in Fig. 1a and converge to results of numerical simulations provided $`ϵ`$ is not too large. For $`\lambda _2`$ numerical evidence indicates the presence of additional non-analytical contributions not captured by the Padé-Borel summation. The results of Fig. 1a allow us to determine the quantity $`\mathrm{\Delta }`$ defined in eq. (2). The result is shown in Fig. 1b. We find that $`\mathrm{\Delta }`$ is maximal for $`ϵ0.21`$ and positive (indicating clustering onto a fractal set) for $`0<ϵ<0.33`$. The red line in Fig. 1d, $`ϵ\kappa \omega ^{1/2}=\text{const}.`$, indicates schematically where $`\mathrm{\Delta }`$ is zero. Above the red line $`\mathrm{\Delta }`$ is always positive, but tends to zero for small $`ϵ`$ as $`\mathrm{\Delta }=10ϵ^2\kappa ^2/\omega `$. In the limit of $`ϵ0`$ the dynamics becomes advective (despite being underdamped): to lowest order in $`ϵ`$ our results coincide with those for purely advective flow LeJ85 .
We now turn to the rate $`J`$ of caustic formation. It is the rate at which $`\delta V(t)=(\delta 𝒓_1(t)\delta 𝒓_2(t))\delta 𝒓_3(t)`$ goes through zero. Since $`\delta 𝒑_\mu `$ typically remain bounded, caustics correspond to instances where the elements of the third column of $`𝐑^{}`$ go to $`\mathrm{}`$ and reappear at $`\mathrm{}`$. The rate at which these events occur is given by the escape rate of the Langevin process (10) to infinity. It is expected Wil05 to have a non-analytic dependence on $`ϵ`$, of the form $`C\mathrm{exp}(S/ϵ^2)`$, as demonstrated in Fig. 1c. In this panel, $`J/\gamma `$ is compared to $`(\lambda _1+\lambda _2+\lambda _3)/\gamma `$. We see that caustics are very rare when $`ϵ1`$, but frequent when $`ϵ`$ is large and they are the only clustering mechanism when $`ϵ>0.33`$.
Finally, we comment on the relation between our results and earlier works (cited above), which suggest that clustering only occurs for $`\omega 1`$ (with the value of $`\kappa `$ unspecified). It must be emphasised that the earlier quantitative theoretical results on clustering are confined to the overdamped limit $`\omega 1`$, where inertial effects are small: for purely advective flow there is no clustering ($`\mathrm{\Delta }=0`$ and $`J=0`$). Inertial effects were incorporated by Elperin and others Elp96 ; Pin99 ; Bal01 , using Maxey’s perturbative correction to the velocity. Their results are valid only for the limit $`\omega 1`$, and are distinct from our series expansions (Clustering by mixing flows): this is most easily seen by calculating corrections to $`\mathrm{\Delta }`$ in this overdamped limit. We find that $`\mathrm{\Delta }\kappa ^2/\omega ^2`$ implying that clustering effects are small in this regime. In the underdamped regime, by contrast, we obtained $`\mathrm{\Delta }\kappa ^2/\omega `$ which can be of order unity.
The results of this letter are summarised schematically in figure 1d. First, at small $`\kappa `$, strong clustering occurs in the region indicated, above the line $`ϵ\kappa \omega ^{1/2}=0.33`$. Second, since the dimension deficit $`\mathrm{\Delta }`$ is positive in this regime, the reasoning of Sommerer and Ott Som93 indicates that the particles cluster on a fractal. Third, as $`ϵ0`$ the dynamics becomes advective. In this limit the dimension deficit $`\mathrm{\Delta }`$ and the rate of caustic formation $`J`$ vanish: particles advected in an incompressible flow remain uniformly distributed. Fourth, when $`ϵ>0.33`$ we find that the dimension deficit $`\mathrm{\Delta }`$ is negative implying that do not lie on a fractal. They are however not homogenously distributed: in this regime particles cluster because they are brought into close contact by caustics.
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# 2005 International Linear Collider Workshop - Stanford, U.S.A. Effects of the 𝑡𝑡̄ threshold in 𝑒⁺𝑒⁻→𝑡𝑡̄𝐻
## I Introduction
It is one of the major tasks of future collider experiments to unravel details of the mechanism of electroweak symmetry breaking (EWSB). In the Standard Model (SM) EWSB is achieved by the Higgs mechanism. The particle masses are generated by the Higgs field vacuum expectation value $`V=(\sqrt{2}G_F)^{1/2}246`$ GeV arising through the Higgs self interactions. The mechanism also predicts that Higgs bosons can be produced in collider experiments. While a Higgs boson with a mass smaller than $`1`$ TeV can be found at the LHC, precise and model-independent measurements of its quantum numbers and couplings can be gained from the $`e^+e^{}`$ Linear Collider. A crucial prediction of the Higgs mechanism is that the Higgs Yukawa coupling to quarks $`\lambda _q`$ is related to the quark masses by $`m_q=\lambda _qV`$. At the $`e^+e^{}`$ Linear Collider the top quark Yukawa coupling can be measured from top quark pair production associated with a Higgs boson, $`e^+e^{}t\overline{t}H`$, since the process is dominated by the amplitudes describing Higgs radiation off the $`t\overline{t}`$ pair. This process is particularly suited for a light Higgs boson since the cross section can then reach the $`1`$-$`2`$ fb level. Assuming an experimental precision at the percent level, QCD and electroweak radiative corrections need to be accounted for in the theoretical predictions. In the approximation that top quarks and the Higgs are stable particles the Born cross section was already determined some time ago in Refs. Borneetth . For the $`𝒪(\alpha _s)`$ QCD one-loop corrections a number of references in various approximations exist Dawson1 ; Dawson2 ; Dittmaier1 . On the other hand, the full set of one-loop electroweak corrections was obtained in Refs. Belanger1 ; Denner1 and also in Refs. You1 . In Ref. Denner1 a detailed analysis of various differential distributions of the cross section $`\sigma (e^+e^{}t\overline{t}H)`$ can be found.
A particularly interesting kinematical phase space region is where the energy of the Higgs boson is large and close to its kinematic endpoint. The $`t\overline{t}`$ pair then becomes collinear and flies opposite to the Higgs direction to balance the large Higgs momentum, see Fig. 1. For large $`E_H`$, on the other hand, the $`t\overline{t}`$ invariant mass $`Q^2`$ approaches $`4m_t^2`$, $`E_H=\frac{1}{2\sqrt{s}}\left(s+m_H^2Q^2\right)`$, and the top quark pair is nonrelativistic in its own center-of-mass (c.m.) system. Because the Higgs is very narrow for a mass below the $`W^+W^{}`$ threshold, strong interactions between the $`t\overline{t}`$ pair and the hadronic Higgs final state can be neglected. Thus, close to the Higgs energy endpoint the $`t\overline{t}`$ QCD dynamics is exclusively governed by the nonrelativistic physics known from the process $`e^+e^{}t\overline{t}`$ in the $`t\overline{t}`$ threshold region at $`\sqrt{s}2m_t`$. In this regime the so-called Coulomb singularities $`(\alpha _s/v)^n`$, with $`v=(14m_t^2/Q^2)^{1/2}`$ being the top quark relative velocity in the $`t\overline{t}`$ c.m. frame, arise and require predictions using an expansion in $`\alpha _s`$ and $`v`$ rather than just a perturbative computation in the number of loops.
This singularity structure is most easily visible in the Higgs energy distribution, $`d\sigma (e^+e^{}t\overline{t}H)/dE_H`$. While the Born distribution approaches zero for $`E_HE_H^{\mathrm{max}}`$, $`d\sigma /dE_Hv`$ Borneetth , the $`𝒪(\alpha _s)`$ fixed-order perturbative corrections are proportional to $`\alpha _s`$ at the endpoint Dawson2 ; Dittmaier1 and the $`𝒪(\alpha _s^2)`$ corrections even diverge like $`\alpha _s^2/v`$. The problem might be avoided by imposing a cut on $`E_H`$ or $`Q^2`$, but this is unnecessary because there exists an elaborate technology being used for the threshold region in the process $`e^+e^{}t\overline{t}`$ TTBARreview that allows for systematic QCD predictions with renormalization group (RG) improvement. Imposing a cut would be also disadvantageous as the nonrelativistic portion of the $`t\overline{t}H`$ phase space increases for smaller c.m. energies which are relevant for a measurement in the first phase of the ILC program. Because the SM top width is quite large, $`\mathrm{\Gamma }_t1.5`$ GeV, the corresponding QCD effective theory computations can be carried out with perturbative methods for all Higgs energies in the endpoint region.
In this talk we present the Higgs energy distribution $`d\sigma /dE_H`$ in the large Higgs energy endpoint region at NLL order in the nonrelativistic expansion using the framework of “velocity” NRQCD (vNRQCD). For details on the conceptual aspects, concerning powercounting, the operator structure of the effective theory action, and renormalization we refer to Refs. LMR ; amis ; amis2 ; HoangStewartultra ; hmst . For a more detailed discussion of the computations for this work see Ref. farrellhoang1 .
## II Review of Effective Theory Ingredients
The effective theory vNRQCD provides a systematic RG-improved description of dynamics of nonrelativistic $`t\overline{t}`$ pairs. The system is characterized, for any energy in the threshold region, by the hierarchy
$`m_tm_tv\text{(three-momentum, “soft” scale)}m_tv^2\text{(kinetic energy, “ultrasoft” scale)}\mathrm{\Lambda }_{\mathrm{QCD}}.`$ (1)
The particle-antiparticle propagation is described by the terms in the effective theory Lagrangian which are bilinear in the top quark and antitop quark fields,
$`(x)`$ $`=`$ $`{\displaystyle \underset{𝒑}{}}\psi _𝒑^{}(x)\left\{i^0{\displaystyle \frac{𝒑^2}{2m_t}}+{\displaystyle \frac{i}{2}}\mathrm{\Gamma }_t\delta m_t\right\}\psi _𝒑(x)+(\psi _𝒑(x)\chi _𝒑(x)),`$ (2)
where the fields $`\psi _𝒑`$ and $`\chi _𝒑`$ destroy top and antitop quarks with soft three-momentum $`𝒑`$ in the $`t\overline{t}`$ c.m. frame and $`\mathrm{\Gamma }_t`$ is the on-shell top quark decay width. The term $`\delta m_t`$ is a residual mass term specific to the top quark mass definition that is being used; for our analysis we employ the 1S mass scheme Hoangupsilon ; HoangTeubnerdist .
Up to NLL order the top-antitop quark pair interacts only through the effective Coulomb potential FischlerBilloire ,
$`\stackrel{~}{V}_c(𝒑,𝒒)`$ $`=`$ $`{\displaystyle \frac{4\pi C_F\alpha _s(m_t\nu )}{𝒌^2}}\left\{\mathrm{\hspace{0.17em}1}+{\displaystyle \frac{\alpha _s(m_t\nu )}{4\pi }}\left[\beta _0\mathrm{ln}\left({\displaystyle \frac{𝒌^2}{m_t^2\nu ^2}}\right)+a_1\right]\right\},`$ (3)
where $`𝒌=𝒑𝒒`$ is the momentum transfer and $`\beta _0=11/3C_A4/3Tn_f`$ is the one-loop QCD beta function, $`a_1=31/9C_A20/9Tn_f`$ the coefficient of the one-loop correction to the effective Coulomb potential, and $`C_A=3,C_F=4/3,T=1/2`$ are SU(3) group theoretical factors. For the number of light flavors we take $`n_f=5`$. The parameter $`\nu `$ is the vNRQCD renormalization scaling parameter used to describe the correlated running of soft and ultrasoft effects in the effective theory. Here, $`\nu =1`$ corresponds to the hard scale at which the effective theory is matched to the full theory, and $`\nu =v_0`$, $`v_0`$ being of the order of the typical $`t\overline{t}`$ relative velocity, is the scale where the matrix elements are computed. The evolution of the Wilson coefficients from the matching scale down to the low-energy scale sums logarithms of the velocity to all orders and is governed by the velocity renormalization group equations LMR .
Top-antitop quark production in the nonrelativistic regime in the LL and NLL approximation in a $`{}_{}{}^{3}S_{1}^{}`$ spin triplet or a $`{}_{}{}^{1}S_{0}^{}`$ spin singlet state is described by the currents
$$J_{1,𝒑}^j=\psi _𝒑^{}\sigma _j(i\sigma _2)\chi _𝒑^{},J_{0,𝒑}=\psi _𝒑^{}(i\sigma _2)\chi _𝒑^{},$$
(4)
where $`c_{1,j}(\nu )`$ and $`c_0(\nu )`$ are the corresponding Wilson coefficients. The currents do not run at LL order but UV-divergences in effective theory two-loop vertex diagrams LMR lead to non-trivial anomalous dimensions at NLL order which result in a scaling of the Wilson coefficients HoangStewartultra ; Pineda1
$$c_{1,j}(\nu )=c_{1,j}(1)\mathrm{exp}\left(f(\nu ,𝐒^\mathrm{𝟐}=2)\right),c_0(\nu )=c_0(1)\mathrm{exp}\left(f(\nu ,𝐒^\mathrm{𝟐}=0)\right),$$
(5)
where $`𝐒^\mathrm{𝟐}`$ is the square of the total $`t\overline{t}`$ spin. We have adopted the convention that the matching conditions at $`\nu =1`$ only account for QCD effects, so at LL order we have $`c_1(1)=c_0(1)=1`$. The NLL order QCD matching conditions relevant for $`e^+e^{}t\overline{t}H`$ in the large Higgs energy endpoint region are discussed in Sec. III.
Through the optical theorem the $`t\overline{t}`$ production rate for a $`t\overline{t}`$ invariant mass $`Q^24m_t^2`$ involves the imaginary part of the time-ordered product of the production and annihilation currents defined in Eqs. (4),
$`𝒜_1^{lk}(Q^2,m_t,\nu )`$ $`=`$ $`i{\displaystyle \underset{𝒑,𝒑^{\mathbf{}}}{}}{\displaystyle d^4xe^{i\widehat{q}x}\mathrm{\hspace{0.17em}0}\left|TJ_{1,𝒑^{}}^l(0)J_{1,𝒑}^k(x)\right|\mathrm{\hspace{0.17em}0}}=\mathrm{\hspace{0.17em}2}N_c\delta ^{lk}G^c(a,v,m_t,\nu ),`$ (6)
$`𝒜_0(Q^2,m_t,\nu )`$ $`=`$ $`i{\displaystyle \underset{𝒑,𝒑^{\mathbf{}}}{}}{\displaystyle d^4xe^{i\widehat{q}x}\mathrm{\hspace{0.17em}0}\left|TJ_{0,𝒑^{}}^{}(0)J_{0,𝒑}(x)\right|\mathrm{\hspace{0.17em}0}}=N_cG^c(a,v,m_t,\nu ),`$ (7)
where $`v=((\sqrt{Q^2}2m_t2\delta m_t+i\mathrm{\Gamma }_t)/m_t)^{\frac{1}{2}}`$ is the c.m. top quark (effective) relative velocity and $`\widehat{q}(\sqrt{Q^2}2m_t,0)`$. The term $`G^c`$ is the zero-distance S-wave Coulomb Green function of the nonrelativistic Schrödinger equation with the potential in Eq. (3). To compute the Green function we use the numerical techniques and codes of the TOPPIC program developed in Ref. Jezabek1 and determine an exact solution of the full NLL Schrödinger equation following the approach of Refs. hmst .
## III Comments on the Computation
The LL vNRQCD result for the Higgs energy distribution, including the QCD effects coming from the Coulomb potential in Eq. (3), the finite top quark lifetime, and the Wilson coefficients of the currents in Eqs. (5), is given by
$`{\displaystyle \frac{d\sigma }{dE_H}}(E_HE_H^{\mathrm{max}})`$ $`=`$ $`{\displaystyle \frac{8N_c\left[(1+x_H4x_t)^24x_H\right]^{1/2}}{s^{3/2}m_t^2}}\left(c_0^2(\nu )F_0^Z+c_1^2(\nu )F_1^{\gamma ,Z}\right)\text{Im}\left[G^c(a,v,m_t,\nu )\right].`$ (8)
For details on the explicit calculation and the definition of the formfactors $`F_i`$ see farrellhoang1 . We note that the NLL ($`𝒪(\alpha _s)`$) matching conditions for the three triplet Wilson coefficients $`c_{1,j}`$ depend on the $`t\overline{t}`$ spin configuration (i.e. on $`j`$) since the kinematic situation for $`t\overline{t}H`$ production in the large Higgs energy endpoint is not invariant under separate rotations of the spin quantization axis. However, for our purposes it is sufficient to define a triplet Wilson coefficient that is averaged over the three spin configurations. Using such an averaged triplet Wilson coefficient the Higgs energy spectrum at NLL order can also be cast in the simple form of Eq. (8).
At NLL order, we need to account for the $`𝒪(\alpha _s^2)`$ contributions to the Coulomb potential in Eq. (3), the NLL running of the coefficients $`c_1`$ and $`c_0`$, and their $`𝒪(\alpha _s)`$ matching conditions at $`\nu =1`$. The latter hard QCD corrections are process specific and cannot be inferred from results obtained in earlier computations for the $`t\overline{t}`$ threshold in $`e^+e^{}t\overline{t}`$.
We have extracted the matching conditions from the codes for the Standard Model amplitude for $`e^+e^{}t\overline{t}H`$ provided in Ref. Denner1 . With the ansatz
$$c_{0,1}(\nu =1,\sqrt{s},m_t,m_H)=1+\frac{C_F\alpha _s(m_t)}{\pi }\delta c_{0,1}(\sqrt{s},m_t,m_H)$$
(9)
we have determined the $`𝒪(\alpha _s)`$ matching conditions numerically by matching the $`𝒪(\alpha _s)`$ vNRQCD prediction at $`\mu =m_t`$ ($`\nu =1`$) to the full theory results close to the large Higgs energy endpoint. The relative uncertainties for this numerical procedure are below 1% for $`\delta c_{0,1}`$.
## IV Numerical Analysis
In Figs. 3 the predictions of the Higgs energy spectrum in the full kinematic range are displayed for two cases. In the large Higgs energy endpoint region we have shown the LL (dashed lines) and NLL (solid lines) results in the nonrelativistic expansion. (See the figure captions for details on choice of parameters.) At LL order the upper (lower) curve corresponds to $`\nu =0.1`$ ($`0.4`$), while at NLL order the upper (lower) curve corresponds to $`\nu =0.2`$ ($`0.1`$). The curves show the typical behavior of the prediction of the nonrelativistic expansion for any choice of parameters. While the LL predictions have a quite large renormalization parameter dependence at the level of several tens of percent, the NLL results are stable. Here, the variation due to change of the renormalization parameter is around 5%. The stabilization with respect to renormalization parameter variations at NLL order arises mainly from the inclusion of the $`𝒪(\alpha _s)`$ QCD corrections to the Coulomb potential, Eq. (3). Moreover, the NLL curves lie considerably lower than the LL ones. This behavior is well known from the predictions for $`e^+e^{}t\overline{t}`$ at threshold TTBARreview ; hmst and originates from the structure of the large negative $`𝒪(\alpha _s)`$ QCD corrections to the Coulomb potential, Eq. (3), and from the sizeable negative QCD corrections to the matching conditions in Eq. (9).
In principle this behavior is a point of concern because it could indicate that the renormalization parameter variation might be an inadequate method to estimate theoretical uncertainties. Fortunately, the top quark mass is sufficiently large such that the regions where the conventional fixed-order expansion (in powers of the strong coupling) and where the nonrelativistic expansion (described by the effective theory) can be applied are expected to overlap. To demonstrate this issue we have also displayed in Figs. 3 the fixed-order (i.e. without summation of Coulomb singularities) predictions at the Born (dashed lines) Dawson1 and the $`𝒪(\alpha _s)`$ level (solid lines) Denner1 . The two $`𝒪(\alpha _s)`$ curves correspond to the renormalization scales $`\mu =\sqrt{s}`$ (lower curves) and $`\mu =|\sqrt{s}v|`$ (upper curves), where $`v`$ is the $`t\overline{t}`$ relative velocity defined below Eq. (7). The latter choice for the fixed-order renormalization scale is motivated by the fact that the relative momentum of the top pair is the scale governing the Coulomb singularities contained in the fixed-order expansion close to the large Higgs energy endpoint. This choice for the fixed-order renormalization scale is therefore the more appropriate one near the Higgs energy endpoint. The results in Figs. 3 demonstrate the overlap between the $`𝒪(\alpha _s)`$ fixed-order prediction and the NLL nonrelativistic one in the region where the $`t\overline{t}`$ relative velocity is approximately 0.2. (The Higgs energy with $`v=0.2`$ is indicated in each panel by the solid vertical line.) The overlap improves for increasing c.m. energies or decreasing Higgs masses. This indicates that in the overlap regions the higher order contributions summed in the nonrelativistic prediction and the higher order relativistic corrections contained in the fixed-order result are both small. For smaller c.m. energies or increasing Higgs masses, on the other hand, the NLL nonrelativistic predictions tend to lie slightly above the $`𝒪(\alpha _s)`$ fixed-order results (for $`\mu =\sqrt{s}v`$) illustrating the impact of the higher order corrections to each type of expansion. The discrepancy, however, remains comparable to the uncertainties estimated from the renormalization parameter variation of the NLL nonrelativistic prediction. We therefore conclude that the renormalization parameter variation of the NLL order nonrelativistic prediction should provide a realistic estimate of the theoretical uncertainties in the large Higgs energy region. The results in Figs. 3 also demonstrate that the region of parameter space where the top quark pair is nonrelativistic increases for smaller c.m. energies (or larger Higgs masses).
Let us now discuss the numerical impact of the nonrelativistic contributions in the large Higgs energy region on the total cross section. In Tab. 1 the importance of the summation of the Coulomb singularities and the logarithms of the top quark velocity is analyzed numerically for various choices of the c.m. energy and the Higgs mass. For all cases the top quark mass $`m_t^{1\mathrm{S}}=180`$ GeV is used and the other parameters are fixed as in Figs. 3. See the caption for details on the various entries. The numbers for $`\sigma (\text{NLL})`$, which are determined from combining the NLL nonrelativistic predictions in the Higgs energy end point for $`|v|<0.2`$ with the fixed-order $`𝒪(\alpha _s)`$ prediction for smaller Higgs energies with $`|v|>0.2`$, represent the currently most complete predictions for the total cross section of the process $`e^+e^{}t\overline{t}H`$ as far as QCD corrections are concerned. For $`\sigma (\text{NLL})`$ we have also given our estimate for the theoretical error. For the fixed-order contribution ($`v>0.2`$) we have estimated the uncertainty by taking the maximum of the shifts obtained from varying $`\mu `$ in the ranges $`[\sqrt{s},2\sqrt{s}]`$, $`[\sqrt{s},\sqrt{s}/2]`$ and $`[\sqrt{s}v,\sqrt{s}]`$; for the nonrelativistic contribution in the end point we have assumed an uncertainty of 5% for all cases. For the numbers displayed in Tab. 1 both uncertainties were added linearly.
The results show that the enhancement of the cross section due to the summations in the large Higgs energy region is particularly important for smaller c.m. energies and larger Higgs masses, when the portion of the phase space where the nonrelativistic expansion has to be applied is large. Here, the higher order summations contained in the nonrelativistic expansion can be comparable to the already sizeable $`𝒪(\alpha _s)`$ fixed-order corrections and enhance the cross section further. This is advantageous for top Yukawa coupling measurements for the lower c.m. energies accessible in the first phase of the ILC experiment. For higher c.m. energies the effect of the nonrelativistic summations of contributions from beyond $`𝒪(\alpha _s)`$ is less pronounced and decreases to the one-percent level for c.m. energies above $`700`$ GeV. For all cases, except for very large c.m. energies around $`1000`$ GeV, however, the shift caused by the terms that are summed up in the nonrelativistic expansion exceeds the theoretical error farrellhoang1 .
In Table 1, in the last column, the ratio of the NLL nonrelativistic cross section and the $`𝒪(\alpha _s)`$ fixed-order cross section (with the approximation $`m_t^{\mathrm{pole}}=m_t^{1\mathrm{S}}`$) for $`|v|<0.2`$ is also shown. Interestingly, the higher order summations lead to correction factors ranging between about $`1.7`$ and $`1.8`$ that are only very weakly dependent of the c.m. energy and the Higgs mass. This fact might prove useful for rough implementations of nonrelativistic $`t\overline{t}`$ effects in other high energy processes.
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# Cohomology of face rings, and torus actions
## 1. Introduction
This article centres on the cohomological aspects of ‘toric topology’, a new and actively developing field on the borders of equivariant topology, combinatorial geometry and commutative algebra. The algebro-geometric counterpart of toric topology, known as ‘toric geometry’ or algebraic geometry of *toric varieties*, is now a well established field in algebraic geometry, which is characterised by its strong links with combinatorial and convex geometry (see the classical survey paper or more modern exposition ). Since the appearance of Davis and Januszkiewicz’s work , where the concept of a *(quasi)toric manifold* was introduced as a topological generalisation of smooth compact toric variety, there has grown an understanding that most phenomena of smooth toric geometry may be modelled in the purely topological situation of smooth manifolds with a nicely behaved torus action.
One of the main results of is that the equivariant cohomology of a toric manifold can be identified with the *face ring* of the quotient simple polytope, or, for more general classes of torus actions, with the face ring of a certain simplicial complex $`K`$. The ordinary cohomology of a quasitoric manifold can also be effectively identified as the quotient of the face ring by a *regular sequence* of degree-two elements, which provides a generalisation to the well-known Danilov–Jurkiewicz theorem of toric geometry. The notion of the face ring of a simplicial complex sits in the heart of Stanley’s ‘Combinatorial commutative algebra’ , linking geometrical and combinatorial problems concerning simplicial complexes with commutative and homological algebra. Our concept of toric topology aims at extending these links and developing new applications by applying the full strength of the apparatus of equivariant topology of torus actions.
The article surveys certain new developments of toric topology related to the cohomology of face rings. Introductory remarks can be found at the beginning of each section and most subsections. A more detailed description of the history of the subject, together with an extensive bibliography, can be found in and its extended Russian version .
The current article represents the work of the algebraic topology and combinatorics group at the Department of Geometry and Topology, Moscow State University, and the author thanks all its members for the collaboration and insight gained from numerous discussions, particularly mentioning Victor Buchstaber, Ilia Baskakov, and Arseny Gadzhikurbanov. The author is also grateful to Nigel Ray for several valuable comments and suggestions that greatly improved this text and his hospitality during the visit to Manchester sponsored by an LMS grant.
## 2. Simplicial complexes and face rings
The notion of the face ring $`𝐤[K]`$ of a simplicial complex $`K`$ is central to the algebraic study of triangulations. In this section we review its main properties, emphasising functoriality with respect to simplicial maps. Then we introduce the bigraded $`\mathrm{Tor}`$-algebra $`\mathrm{Tor}_{𝐤[v_1,\mathrm{},v_m]}(𝐤[K],𝐤)`$ through a finite free resolution of $`𝐤[K]`$ as a module over the polynomial ring. The corresponding bigraded Betti numbers are important combinatorial invariants of $`K`$.
### 2.1. Definition and main properties
Let $`K=K^{n1}`$ be an arbitrary $`(n1)`$-dimensional simplicial complex on an $`m`$-element vertex set $`V`$, which we usually identify with the set of ordinals $`[m]=\{1,\mathrm{},m\}`$. Those subsets $`\sigma V`$ belonging to $`K`$ are referred to as *simplices*; we also use the notation $`\sigma K`$. We count the empty set $`\mathrm{}`$ as a simplex of $`K`$. When it is necessary to distinguish between combinatorial and geometrical objects, we denote by $`|K|`$ a *geometrical realisation* of $`K`$, which is a triangulated topological space.
Choose a ground commutative ring $`𝐤`$ with unit (we are mostly interested in the cases $`𝐤=,`$ or finite field). Let $`𝐤[v_1,\mathrm{},v_m]`$ be the graded polynomial algebra over $`𝐤`$ with $`\mathrm{deg}v_i=2`$. For an arbitrary subset $`\omega =\{i_1,\mathrm{},i_k\}[m]`$, denote by $`v_\omega `$ the square-free monomial $`v_{i_1}\mathrm{}v_{i_k}`$.
The *face ring* (or *Stanley–Reisner algebra*) of $`K`$ is the quotient ring
$$𝐤[K]=𝐤[v_1,\mathrm{},v_m]/_K,$$
where $`_K`$ is the homogeneous ideal generated by all monomials $`v_\sigma `$ such that $`\sigma `$ is not a simplex of $`K`$. The ideal $`_K`$ is called the *Stanley–Reisner ideal* of $`K`$.
###### Example 2.1.
Let $`K`$ be a 2-dimensional simplicial complex shown on Figure 1. Then
$$𝐤[K]=𝐤[v_1,\mathrm{},v_5]/(v_1v_5,v_3v_4,v_1v_2v_3,v_2v_4v_5).$$
Despite its simple construction, the face ring appears to be a very powerful tool allowing us to translate the combinatorial properties of different particular classes of simplicial complexes into the language of commutative algebra. The resulting field of ‘Combinatorial commutative algebra’, whose foundations were laid by Stanley in his monograph , has attracted a lot of interest from both combinatorialists and commutative algebraists.
Let $`K_1`$ and $`K_2`$ be two simplicial complexes on the vertex sets $`[m_1]`$ and $`[m_2]`$ respectively. A set map $`\phi :[m_1][m_2]`$ is called a *simplicial map* between $`K_1`$ and $`K_2`$ if $`\phi (\sigma )K_2`$ for any $`\sigma K_1`$; we often identify such $`\phi `$ with its restriction to $`K_1`$ (regarded as a collection of subsets of $`[m_1]`$), and use the notation $`\phi :K_1K_2`$.
###### Proposition 2.2.
Let $`\phi :K_1K_2`$ be a simplicial map. Define a map $`\phi ^{}:𝐤[w_1,\mathrm{},w_{m_2}]𝐤[v_1,\mathrm{},v_{m_1}]`$ by
$$\phi ^{}(w_j):=\underset{i\phi ^1(j)}{}v_i.$$
Then $`\phi ^{}`$ induces a homomorphism $`𝐤[K_2]𝐤[K_1]`$, which we will also denote by $`\phi ^{}`$.
###### Proof.
We have to check that $`\phi ^{}(_{K_2})_{K_1}`$. Suppose $`\tau =\{j_1,\mathrm{},j_s\}[m_2]`$ is not a simplex of $`K_2`$. Then
(2.1)
$$\phi ^{}(w_{j_1}\mathrm{}w_{j_s})=\underset{i_1\phi ^1(j_1),\mathrm{},i_s\phi ^1(j_s)}{}v_{i_1}\mathrm{}v_{i_s}.$$
We claim that $`\sigma =\{i_1,\mathrm{},i_s\}`$ is not a simplex of $`K_1`$ for any monomial $`v_{i_1}\mathrm{}v_{i_s}`$ in the right hand side of the above identity. Indeed, if $`\sigma K_1`$, then $`\phi (\sigma )=\tau K_2`$ by the definition of simplicial map, which leads to a contradiction. Hence, the right hand side of (2.1) is in $`_{K_1}`$. ∎
### 2.2. Cohen–Macaulay rings and complexes
Cohen–Macaulay rings and modules play an important role in homological commutative algebra and algebraic geometry. A standard reference for the subject is , where the reader may find proofs of the basic facts about Cohen–Macaulay rings and regular sequences mentioned in this subsection. In the case of simplicial complexes, the Cohen–Macaulay property of the corresponding face rings leads to important combinatorial and topological consequences.
Let $`A=_{i0}A^i`$ be a finitely-generated commutative graded algebra over $`𝐤`$. We assume that $`A`$ is connected ($`A^0=𝐤`$) and has only even-degree graded components, so that we do not need to distinguish between graded and non-graded commutativity. We denote by $`A_+`$ the positive-degree part of $`A`$ and by $`(A_+)`$ the set of homogeneous elements in $`A_+`$.
A sequence $`t_1,\mathrm{},t_n`$ of algebraically independent homogeneous elements of $`A`$ is called an *hsop* (homogeneous system of parameters) if $`A`$ is a finitely-generated $`𝐤[t_1,\mathrm{},t_n]`$-module (equivalently, $`A/(t_1,\mathrm{},t_n)`$ has finite dimension as a $`𝐤`$-vector space).
###### Lemma 2.3 (Nöther normalisation lemma).
Any finitely-generated graded algebra $`A`$ over a field $`𝐤`$ admits an hsop. If $`𝐤`$ has characteristic zero and $`A`$ is generated by degree-two elements, then a degree-two hsop can be chosen.
A degree-two hsop is called an *lsop* (linear system of parameters).
A sequence $`\text{t}=t_1,\mathrm{},t_k`$ of elements of $`(A_+)`$ is called a *regular sequence* if $`t_{i+1}`$ is not a zero divisor in $`A/(t_1,\mathrm{},t_i)`$ for $`0i<k`$. A regular sequence consists of algebraically independent elements, so it generates a polynomial subring in $`A`$. It can be shown that t is a regular sequence if and only if $`A`$ is a *free* $`𝐤[t_1,\mathrm{},t_k]`$-module.
An algebra $`A`$ is called *Cohen–Macaulay* if it admits a regular hsop t. It follows that $`A`$ is Cohen–Macaulay if and only if it is a free and finitely generated module over its polynomial subring. If $`𝐤`$ is a field of zero characteristic and $`A`$ is generated by degree-two elements, then one can choose t to be an lsop. A simplicial complex $`K`$ is called *Cohen–Macaulay* (over $`𝐤`$) if its face ring $`𝐤[K]`$ is Cohen–Macaulay.
###### Example 2.4.
Let $`K=\mathrm{\Delta }^2`$ be the boundary of a 2-simplex. Then
$$𝐤[K]=𝐤[v_1,v_2,v_3]/(v_1v_2v_3).$$
The elements $`v_1,v_2𝐤[K]`$ are algebraically independent, but do not form an hsop, since $`𝐤[K]/(v_1,v_2)𝐤[v_3]`$ is not finite-dimensional as a $`𝐤`$-space. On the other hand, the elements $`t_1=v_1v_3`$, $`t_2=v_2v_3`$ of $`𝐤[K]`$ form an hsop, since $`𝐤[K]/(t_1,t_2)𝐤[t]/t^3`$. It is easy to see that $`𝐤[K]`$ is a free $`𝐤[t_1,t_2]`$-module with one 0-dimensional generator 1, one 1-dimensional generator $`v_1`$, and one 2-dimensional generator $`v_1^2`$. Thus, $`𝐤[K]`$ is Cohen–Macaulay and $`(t_1,t_2)`$ is a regular sequence.
For an arbitrary simplex $`\sigma K`$ define its *link* and *star* as the subcomplexes
$`link_K\sigma `$ $`=\{\tau K:\sigma \tau K,\sigma \tau =\mathrm{}\};`$
$`star_K\sigma `$ $`=\{\tau K:\sigma \tau K\}.`$
If $`vK`$ is a vertex, then $`star_Kv`$ is the subcomplex consisting of all simplices of $`K`$ containing $`v`$, and all their subsimplices. Note also that $`star_Kv`$ is the cone over $`link_Kv`$.
The following fundamental theorem characterises Cohen–Macaulay complexes combinatorially.
###### Theorem 2.5 (Reisner).
A simplicial complex $`K`$ is Cohen–Macaulay over $`𝐤`$ if and only if for any simplex $`\sigma K`$ (including $`\sigma =\mathrm{}`$) and $`i<dim(link_K\sigma )`$, it holds that $`\stackrel{~}{H}_i(link_K\sigma ;𝐤)=0`$.
Using standard techniques of $`PL`$ topology the previous theorem may be reformulated in purely topological terms.
###### Proposition 2.6 (Munkres).
$`K^{n1}`$ is Cohen–Macaulay over $`𝐤`$ if and only if for an arbitrary point $`x|K|`$, it holds that
$$\stackrel{~}{H}_i(|K|;𝐤)=H_i(|K|,|K|\backslash x;𝐤)=0\text{ for }i<n1.$$
Thus any triangulation of a sphere is a Cohen–Macaulay complex.
### 2.3. Resolutions and $`\mathrm{Tor}`$-algebras
Let $`M`$ be a finitely-generated graded $`𝐤[v_1,\mathrm{},v_m]`$-module. A *free resolution* of $`M`$ is an exact sequence
(2.2)
$$\begin{array}{ccccccccccc}\mathrm{}& \stackrel{d}{}& R^i& \stackrel{d}{}& \mathrm{}& \stackrel{d}{}& R^1& \stackrel{d}{}& R^0& & M0,\end{array}$$
where the $`R^i`$ are finitely-generated graded free $`𝐤[v_1,\mathrm{},v_m]`$-modules and the maps $`d`$ are degree-preserving. By the Hilbert syzygy theorem, there is a free resolution of $`M`$ with $`R^i=0`$ for $`i>m`$. A resolution (2.2) determines a bigraded differential $`𝐤`$-module $`[R,d]`$, where $`R=R^{i,j}`$, $`R^{i,j}:=(R^i)^j`$ and $`d:R^{i,j}R^{i+1,j}`$. The bigraded cohomology module $`H[R,d]`$ has $`H^{i,k}[R,d]=0`$ for $`i>0`$ and $`H^{0,k}[R,d]=M^k`$. Let $`[M,0]`$ be the bigraded module with $`M^{i,k}=0`$ for $`i>0`$, $`M^{0,k}=M^k`$, and zero differential. Then the resolution (2.2) determines a bigraded map $`[R,d][M,0]`$ inducing an isomorphism in cohomology.
Let $`N`$ be another module; then applying the functor $`_{𝐤[v_1,\mathrm{},v_m]}N`$ to a resolution $`[R,d]`$ we get a homomorphism of differential modules
$$[R_{𝐤[v_1,\mathrm{},v_m]}N,d][M_{𝐤[v_1,\mathrm{},v_m]}N,0],$$
which in general does not induce an isomorphism in cohomology. The $`(i)`$th cohomology module of the cochain complex
$$\begin{array}{ccccccccc}\mathrm{}& & R^i_{𝐤[v_1,\mathrm{},v_m]}N& & \mathrm{}& & R^0_{𝐤[v_1,\mathrm{},v_m]}N& & 0\end{array}$$
is denoted by $`\mathrm{Tor}_{𝐤[v_1,\mathrm{},v_m]}^i(M,N)`$. Thus,
$$\mathrm{Tor}_{𝐤[v_1,\mathrm{},v_m]}^i(M,N):=\frac{Ker[d:R^i_{𝐤[v_1,\mathrm{},v_m]}NR^{i+1}_{𝐤[v_1,\mathrm{},v_m]}N]}{d(R^{i1}_{𝐤[v_1,\mathrm{},v_m]}N)}.$$
Since all the $`R^i`$ and $`N`$ are graded modules, we actually have a *bigraded* $`𝐤`$-module
$$\mathrm{Tor}_{𝐤[v_1,\mathrm{},v_m]}(M,N)=\underset{i,j}{}\mathrm{Tor}_{𝐤[v_1,\mathrm{},v_m]}^{i,j}(M,N).$$
The following properties of $`\mathrm{Tor}_{𝐤[v_1,\mathrm{},v_m]}^i(M,N)`$ are well known.
###### Proposition 2.7.
(a) the module $`\mathrm{Tor}_{𝐤[v_1,\mathrm{},v_m]}^i(M,N)`$ does not depend on a choice of resolution in (2.2);
(b) $`\mathrm{Tor}_{𝐤[v_1,\mathrm{},v_m]}^i(,N)`$ and $`\mathrm{Tor}_{𝐤[v_1,\mathrm{},v_m]}^i(M,)`$ are covariant functors;
(c) $`\mathrm{Tor}_{𝐤[v_1,\mathrm{},v_m]}^0(M,N)M_{𝐤[v_1,\mathrm{},v_m]}N`$;
(d) $`\mathrm{Tor}_{𝐤[v_1,\mathrm{},v_m]}^i(M,N)\mathrm{Tor}_{𝐤[v_1,\mathrm{},v_m]}^i(N,M)`$.
Now put $`M=𝐤[K]`$ and $`N=𝐤`$. Since $`\mathrm{deg}v_i=2`$, we have
$$\mathrm{Tor}_{𝐤[v_1,\mathrm{},v_m]}(𝐤[K],𝐤)=\underset{i,j=0}{\overset{m}{}}\mathrm{Tor}_{𝐤[v_1,\mathrm{},v_m]}^{i,2j}(𝐤[K],𝐤)$$
Define the *bigraded Betti numbers* of $`𝐤[K]`$ by
(2.3)
$$\beta ^{i,2j}\left(𝐤[K]\right):=dim_𝐤\mathrm{Tor}_{𝐤[v_1,\mathrm{},v_m]}^{i,2j}(𝐤[K],𝐤),0i,jm.$$
We also set
$$\beta ^i(𝐤[K])=dim_𝐤\mathrm{Tor}_{𝐤[v_1,\mathrm{},v_m]}^i(𝐤[K],𝐤)=\underset{j}{}\beta ^{i,2j}(𝐤[K]).$$
###### Example 2.8.
Let $`K`$ be the boundary of a square. Then
$$𝐤[K]𝐤[v_1,\mathrm{},v_4]/(v_1v_3,v_2v_4).$$
Let us construct a resolution of $`𝐤[K]`$ and calculate the corresponding bigraded Betti numbers. The module $`R^0`$ has one generator 1 (of degree 0), and the map $`R^0𝐤[K]`$ is the quotient projection. Its kernel is the ideal $`_K`$, generated by two monomials $`v_1v_3`$ and $`v_2v_4`$. Take $`R^1`$ to be a free module on two 4-dimensional generators, denoted $`v_{13}`$ and $`v_{24}`$, and define $`d:R^1R^0`$ by sending $`v_{13}`$ to $`v_1v_3`$ and $`v_{24}`$ to $`v_2v_4`$. Its kernel is generated by one element $`v_2v_4v_{13}v_1v_3v_{24}`$. Hence, $`R^2`$ has one generator of degree 8, say $`a`$, and the map $`d:R^2R^1`$ is injective and sends $`a`$ to $`v_2v_4v_{13}v_1v_3v_{24}`$. Thus, we have a resolution
$$\begin{array}{ccccccccccc}0& & R^2& & R^1& & R^0& & M& & 0\end{array}$$
where $`rankR^0=\beta ^{0,0}(𝐤[K])=1`$$`rankR^1=\beta ^{1,4}=2`$ and $`rankR^2=\beta ^{2,8}=1`$.
The Betti numbers $`\beta ^{i,2j}(𝐤[K])`$ are important combinatorial invariants of the simplicial complex $`K`$. The following result expresses them in terms of homology groups of subcomplexes of $`K`$.
Given a subset $`\omega [m]`$, we may restrict $`K`$ to $`\omega `$ and consider the *full subcomplex* $`K_\omega =\{\sigma K:\sigma \omega \}`$.
###### Theorem 2.9 (Hochster).
We have
$$\beta ^{i,2j}\left(𝐤[K]\right)=\underset{\omega [m]:|\omega |=j}{}dim_𝐤\stackrel{~}{H}^{ji1}(K_\omega ;𝐤),$$
where $`\stackrel{~}{H}^{}()`$ denotes the reduced cohomology groups and we assume that $`\stackrel{~}{H}^1(\mathrm{})=𝐤`$.
Hochster’s original proof of this theorem uses rather complicated combinatorial and commutative algebra techniques. Later in subsection 5.1 we give a topological interpretation of the numbers $`\beta ^{i,2j}(𝐤[K])`$ as the bigraded Betti numbers of a topological space, and prove a generalisation of Hochster’s theorem.
###### Example 2.10 (Koszul resolution).
Let $`M=𝐤`$ with the $`𝐤[v_1,\mathrm{},v_m]`$-module structure defined via the map $`𝐤[v_1,\mathrm{},v_m]𝐤`$ sending each $`v_i`$ to 0. Let $`\mathrm{\Lambda }[u_1,\mathrm{},u_m]`$ denote the exterior $`𝐤`$-algebra on $`m`$ generators. The tensor product $`R=\mathrm{\Lambda }[u_1,\mathrm{},u_m]𝐤[v_1,\mathrm{},v_m]`$ (here and below we use $``$ for $`_𝐤`$) may be turned to a differential bigraded algebra by setting
$$bidegu_i=(1,2),bidegv_i=(0,2),$$
(2.4)
$$du_i=v_i,dv_i=0,$$
and requiring $`d`$ to be a derivation of algebras. An explicit construction of a cochain homotopy shows that $`H^i[R,d]=0`$ for $`i>0`$ and $`H^0[R,d]=𝐤`$. Since $`\mathrm{\Lambda }[u_1,\mathrm{},u_m]𝐤[v_1,\mathrm{},v_m]`$ is a free $`𝐤[v_1,\mathrm{},v_m]`$-module, it determines a free resolution of $`𝐤`$. It is known as the *Koszul resolution* and its expanded form (2.2) is as follows:
$$\begin{array}{c}0\mathrm{\Lambda }^m[u_1,\mathrm{},u_m]𝐤[v_1,\mathrm{},v_m]\mathrm{}\hfill \\ \hfill \mathrm{\Lambda }^1[u_1,\mathrm{},u_m]𝐤[v_1,\mathrm{},v_m]𝐤[v_1,\mathrm{},v_m]𝐤0\end{array}$$
where $`\mathrm{\Lambda }^i[u_1,\mathrm{},u_m]`$ is the subspace of $`\mathrm{\Lambda }[u_1,\mathrm{},u_m]`$ spanned by monomials of length $`i`$.
Now let us consider the differential bigraded algebra $`[\mathrm{\Lambda }[u_1,\mathrm{},u_m]𝐤[K],d]`$ with $`d`$ defined as in (2.4).
###### Lemma 2.11.
There is an isomorphism of bigraded modules:
$$\mathrm{Tor}_{𝐤[v_1,\mathrm{},v_m]}(𝐤[K],𝐤)H[\mathrm{\Lambda }[u_1,\mathrm{},u_m]𝐤[K],d]$$
which endows $`\mathrm{Tor}_{𝐤[v_1,\mathrm{},v_m]}(𝐤[K],𝐤)`$ with a bigraded algebra structure in a canonical way.
###### Proof.
Using the Koszul resolution in the definition of $`\mathrm{Tor}`$, we calculate
$`\mathrm{Tor}_{𝐤[v_1,\mathrm{},v_m]}(𝐤[K],𝐤)`$ $`\mathrm{Tor}_{𝐤[v_1,\mathrm{},v_m]}(𝐤,𝐤[K])`$
$`=H\left[\mathrm{\Lambda }[u_1,\mathrm{},u_m]𝐤[v_1,\mathrm{},v_m]_{𝐤[v_1,\mathrm{},v_m]}𝐤[K]\right]`$
$`H\left[\mathrm{\Lambda }[u_1,\mathrm{},u_m]𝐤[K]\right].`$
The cohomology in the right hand side is a bigraded algebra, providing an algebra structure for $`\mathrm{Tor}_{𝐤[v_1,\mathrm{},v_m]}(𝐤[K],𝐤)`$. ∎
The bigraded algebra $`\mathrm{Tor}_{𝐤[v_1,\mathrm{},v_m]}(𝐤[K],𝐤)`$ is called the *$`\mathrm{Tor}`$-algebra* of the simplicial complex $`K`$.
###### Lemma 2.12.
A simplicial map $`\phi :K_1K_2`$ between two simplicial complexes on the vertex sets $`[m_1]`$ and $`[m_2]`$ respectively induces a homomorphism
(2.5)
$$\phi _t^{}:\mathrm{Tor}_{𝐤[w_1,\mathrm{},w_{m_2}]}(𝐤[K_2],𝐤)\mathrm{Tor}_{𝐤[v_1,\mathrm{},v_{m_1}]}(𝐤[K_1],𝐤)$$
of the corresponding $`\mathrm{Tor}`$-algebras.
###### Proof.
This follows directly from Propositions 2.2 and 2.7 (b). ∎
## 3. Toric spaces
Moment-angle complexes provide a functor $`K𝒵_K`$ from the category of simplicial complexes and simplicial maps to the category of spaces with torus action and equivariant maps. This functor allows us to use the techniques of equivariant topology in the study of combinatorics of simplicial complexes and commutative algebra of their face rings; in a way, it breathes a geometrical life into Stanley’s ‘combinatorial commutative algebra’. In particular, the calculation of the cohomology of $`𝒵_K`$ opens a way to a topological treatment of homological invariants of face rings.
The space $`𝒵_K`$ was introduced for arbitrary finite simplicial complex $`K`$ by Davis and Januszkiewicz as a technical tool in their study of (quasi)toric manifolds, a topological generalisation of smooth algebraic toric varieties. Later this space turned out to be of great independent interest. For the subsequent study of $`𝒵_K`$, its place within ‘toric topology’, and connections with combinatorial problems we refer to and its extended Russian version . Here we review the most important aspects of this study related to the cohomology of face rings.
### 3.1. Moment-angle complexes
The $`m`$*-torus* $`T^m`$ is a product of $`m`$ circles; we usually regard it as embedded in $`^m`$ in the standard way:
$$T^m=\{(z_1,\mathrm{},z_m)^m:|z_i|=1,i=1,\mathrm{},m\}.$$
It is contained in the *unit polydisk*
$$(D^2)^m=\{(z_1,\mathrm{},z_m)^m:|z_i|1,i=1,\mathrm{},m\}.$$
For an arbitrary subset $`\omega V`$, define
$$B_\omega :=\{(z_1,\mathrm{},z_m)(D^2)^m:|z_i|=1\text{ for }i\omega \}.$$
The subspace $`B_\omega `$ is homeomorphic to $`(D^2)^{|\omega |}\times T^{m|\omega |}`$.
Given a simplicial complex $`K`$ on $`[m]=\{1,\mathrm{},m\}`$, we define the *moment-angle complex* $`𝒵_K`$ by
(3.1)
$$𝒵_K:=\underset{\sigma K}{}B_\sigma (D^2)^m.$$
The torus $`T^m`$ acts on $`(D^2)^m`$ coordinatewise and each subspace $`B_\omega `$ is invariant under this action. Therefore, the space $`𝒵_K`$ inherits a torus action. The quotient $`(D^2)^m/T^m`$ can be identified with the *unit $`m`$-cube*:
$$I^m:=\{(y_1,\mathrm{},y_m)^m:0y_i1,i=1,\mathrm{},m\}.$$
The quotient $`B_\omega /T^m`$ is then the following $`|\omega |`$-dimensional face of $`I^m`$:
$$C_\omega :=\{(y_1,\mathrm{},y_m)I^m:y_i=1\text{ if }i\omega \}.$$
Thus the whole quotient $`𝒵_K/T^m`$ is identified with a certain cubical subcomplex in $`I^m`$, which we denote by $`cc(K)`$.
###### Lemma 3.1.
The cubical complex $`cc(K)`$ is $`PL`$-homeomorphic to $`coneK`$.
###### Proof.
Let $`K^{}`$ denote the barycentric subdivision of $`K`$ (the vertices of $`K^{}`$ correspond to non-empty simplices $`\sigma `$ of $`K`$). We define a $`PL`$ embedding $`i_c:coneK^{}I^m`$ by mapping each vertex $`\sigma `$ to the vertex $`(\epsilon _1,\mathrm{},\epsilon _m)I^m`$ where $`\epsilon _i=0`$ if $`i\sigma `$ and $`\epsilon _i=1`$ otherwise, the cone vertex to $`(1,\mathrm{},1)I^m`$, and then extending linearly on the simplices of $`coneK^{}`$. The barycentric subdivision of a face $`\sigma K`$ is a subcomplex in $`K^{}`$, which we denote $`K^{}|_\sigma `$. Under the map $`i_c`$ the subcomplex $`coneK^{}|_\sigma `$ maps onto the face $`C_\sigma I^m`$. Thus the whole complex $`coneK^{}`$ maps homeomorphically onto $`cc(K)`$, which concludes the proof. ∎
It follows that the moment-angle complex $`𝒵_K`$ can be defined by the pullback diagram
$$\begin{array}{ccc}𝒵_K& & (D^2)^m\\ & & \rho & & \\ coneK^{}& \stackrel{i_c}{}& I^m\end{array}$$
where $`\rho `$ is the projection onto the orbit space.
###### Example 3.2.
The embedding $`i_c`$ for two simple cases when $`K`$ is a three point complex and the boundary of a triangle is shown on Figure 2. If $`K=\mathrm{\Delta }^{m1}`$ is the whole simplex on $`m`$ vertices, then $`cc(K)`$ is the whole cube $`I^m`$, and the above constructed $`PL`$-homeomorphism between $`cone(\mathrm{\Delta }^{m1})^{}`$ and $`I^m`$ defines the *standard triangulation* of $`I^m`$.
The next lemma shows that the space $`𝒵_K`$ is particularly nice for certain geometrically important classes of triangulations.
###### Lemma 3.3.
Suppose that $`K`$ is a triangulation of an $`(n1)`$-dimensional sphere. Then $`𝒵_K`$ is a closed $`(m+n)`$-dimensional manifold.
In general, if $`K`$ is a triangulated manifold then $`𝒵_K\rho ^1(1,\mathrm{},1)`$ is a noncompact manifold, where $`(1,\mathrm{},1)I^m`$ is the cone vertex and $`\rho ^1(1,\mathrm{},1)T^m`$.
###### Proof.
We only prove the first statement here; the proof of the second is similar and can be found in . Each vertex $`v_i`$ of $`K`$ corresponds to a vertex of the barycentric subdivision $`K^{}`$, which we continue to denote $`v_i`$. Let $`star_K^{}v_i`$ be the star of $`v_i`$ in $`K^{}`$, that is, the subcomplex consisting of all simplices of $`K^{}`$ containing $`v_i`$, and all their subsimplices. The space $`coneK^{}`$ has a canonical *face structure* whose facets (codimension-one faces) are
(3.2)
$$F_i:=star_K^{}v_i,i=1,\mathrm{},m,$$
and whose $`i`$-faces are non-empty intersections of $`i`$-tuples of facets. In particular, the vertices (0-faces) in this face structure are the barycentres of $`(n1)`$-dimensional simplices of $`K`$.
For every such barycentre $`b`$ we denote by $`U_b`$ the subset of $`coneK^{}`$ obtained by removing all faces not containing $`b`$. Since $`K`$ is a triangulation of a sphere, $`coneK^{}`$ is an $`n`$-ball, hence each $`U_b`$ is homeomorphic to an open subset in $`I^n`$ via a homeomorphism preserving the dimension of faces. Since each point of $`coneK^{}`$ is contained in some $`U_b`$, this displays $`coneK^{}`$ as a *manifold with corners*. Having identified $`coneK^{}`$ with $`cc(K)`$ and further $`cc(K)`$ with $`𝒵_K/T^m`$, we see that every point in $`𝒵_K`$ lies in a neighbourhood homeomorphic to an open subset in $`(D^2)^n\times T^{mn}`$ and thus in $`^{m+n}`$. ∎
A particularly important class of examples of sphere triangulations arise from boundary triangulations of convex polytopes. Suppose $`P`$ is a *simple* $`n`$-dimensional convex polytope, i.e. one where every vertex is contained in exactly $`n`$ facets. Then the dual (or *polar*) polytope is *simplicial*, and we denote its boundary complex by $`K_P`$. $`K_P`$ is then a triangulation of an $`(n1)`$-sphere. The faces of $`coneK_P^{}`$ introduced in the previous proof coincide with those of $`P`$.
###### Example 3.4.
Let $`K=\mathrm{\Delta }^{m1}`$. Then $`𝒵_K=((D^2)^m)S^{2m1}`$. In particular, for $`m=2`$ from (3.1) we get the familiar decomposition
$$S^3=D^2\times S^1S^1\times D^2D^2\times D^2$$
of a 3-sphere into a union of two solid tori.
Using faces (3.2) we can identify the isotropy subgroups of the $`T^m`$-action on $`𝒵_K`$. Namely, the isotropy subgroup of a point $`x`$ in the orbit space $`coneK^{}`$ is the coordinate subtorus
$$T(x)=\{(z_1,\mathrm{},z_m)T^m:z_i=1\text{ if }xF_i\}.$$
In particular, the action is free over the interior (that is, near the cone point) of $`coneK^{}`$.
It follows that the moment-angle complex can be identified with the quotient
$$𝒵_K=(T^m\times |coneK^{}|)/,$$
where $`(t_1,x)(t_2,y)`$ if and only if $`x=y`$ and $`t_1t_2^1T(x)`$. In the case when $`K`$ is the dual triangulation of a simple polytope $`P^n`$ we may write $`(T^m\times P^n)/`$ instead. The latter $`T^m`$-manifold is the one introduced by Davis and Januszkiewicz , which thereby coincides with our moment-angle complex.
### 3.2. Homotopy fibre construction
The classifying space for the circle $`S^1`$ can be identified with the infinite-dimensional projective space $`P^{\mathrm{}}`$. The classifying space $`BT^m`$ of the $`m`$-torus is a product of $`m`$ copies of $`P^{\mathrm{}}`$. The cohomology of $`BT^m`$ is the polynomial ring $`[v_1,\mathrm{},v_m]`$, $`\mathrm{deg}v_i=2`$ (the cohomology is taken with integer coefficients, unless another coefficient ring is explicitly specified). The total space $`ET^m`$ of the universal principal $`T^m`$-bundle over $`BT^m`$ can be identified with the product of $`m`$ infinite-dimensional spheres.
In Davis and Januszkiewicz considered the *homotopy quotient* of $`𝒵_K`$ by the $`T^m`$-action (also known as the *Borel construction*). We refer to it as the *Davis–Januszkiewicz space*:
$$\text{DJ}(K):=ET^m\times _{T^m}𝒵_K=ET^m\times 𝒵_K/,$$
where $`(e,z)(et^1,tz)`$. There is a a fibration $`p:\text{DJ}(K)BT^m`$ with fibre $`𝒵_K`$. The cohomology of the Borel construction of a $`T^m`$-space $`X`$ is called the *equivariant cohomology* and denoted by $`H_{T^m}^{}(X)`$.
A theorem of states that the cohomology ring of $`\text{DJ}(K)`$ (or the equivariant cohomology of $`𝒵_K`$) is isomorphic to $`[K]`$. This result can be clarified by an alternative construction of $`\text{DJ}(K)`$ , which we review below.
The space $`BT^m`$ has the canonical cell decomposition in which each factor $`P^{\mathrm{}}`$ has one cell in every even dimension. Given a subset $`\omega [m]`$, define the subproduct
$$BT^\omega :=\{(x_1,\mathrm{},x_m)BT^m:x_i=\text{ if }i\omega \}$$
where $``$ is the basepoint (zero-cell) of $`P^{\mathrm{}}`$. Now for a simplicial complex $`K`$ on $`[m]`$ define the following cellular subcomplex:
(3.3)
$$BT^K:=\underset{\sigma K}{}BT^\sigma BT^m.$$
###### Proposition 3.5.
The cohomology of $`BT^K`$ is isomorphic to the Stanley–Reisner ring $`[K]`$. Moreover, the inclusion of cellular complexes $`i:BT^KBT^m`$ induces the quotient epimorphism
$$i^{}:[v_1,\mathrm{},v_m][K]=[v_1,\mathrm{},v_m]/_K$$
in the cohomology.
###### Proof.
Let $`B_i^{2k}`$ denote the $`2k`$-dimensional cell in the $`i`$th factor of $`BT^m`$, and $`C^{}(BT^m)`$ the cellular cochain module. A monomial $`v_{i_1}^{k_1}\mathrm{}v_{i_p}^{k_p}`$ represents the cellular cochain $`(B_{i_1}^{2k_1}\mathrm{}B_{i_p}^{2k_p})^{}`$ in $`C^{}(BT^m)`$. Under the cochain homomorphism induced by the inclusion $`BT^KBT^m`$ the cochain $`(B_{i_1}^{2k_1}\mathrm{}B_{i_p}^{2k_p})^{}`$ maps identically if $`\{i_1,\mathrm{},i_p\}K`$ and to zero otherwise, whence the statement follows. ∎
###### Theorem 3.6.
There is a deformation retraction $`\text{DJ}(K)BT^K`$ such that the diagram
$$\begin{array}{ccc}\text{DJ}(K)& \stackrel{p}{}& BT^m\\ & & & & \\ BT^K& \stackrel{i}{}& BT^m\end{array}$$
is commutative.
###### Proof.
We have $`𝒵_K=_{\sigma K}B_\sigma `$, and each $`B_\sigma `$ is $`T^m`$-invariant. Hence, there is the corresponding decomposition of the Borel construction:
$$\text{DJ}(K)=ET^m\times _{T^m}𝒵_K=\underset{\sigma K}{}ET^m\times _{T^m}B_\sigma .$$
Suppose $`|\sigma |=s`$. Then $`B_\sigma (D^2)^s\times T^{ms}`$, so we have
$$ET^m\times _{T^m}B_\sigma (ET^s\times _{T^s}(D^2)^s)\times ET^{ms}.$$
The space $`ET^s\times _{T^s}(D^2)^s`$ is the total space of a $`(D^2)^s`$-bundle over $`BT^s`$, and $`ET^{ms}`$ is contractible. It follows that there is a deformation retraction $`ET^m\times _{T^m}B_\sigma BT^\sigma `$. These homotopy equivalences corresponding to different simplices fit together to yield the required homotopy equivalence between $`p:\text{DJ}(K)BT^m`$ and $`i:BT^KBT^m`$. ∎
###### Corollary 3.7.
The space $`𝒵_K`$ is the homotopy fibre of the cellular inclusion $`i:BT^KBT^m`$. Hence there are ring isomorphisms
$$H^{}(\text{DJ}(K))=H_{T^m}^{}(𝒵_K)[K].$$
In view of the last two statements we shall also use the notation $`\text{DJ}(K)`$ for $`BT^K`$, and refer to the whole class of spaces homotopy equivalent to $`\text{DJ}(K)`$ as the *Davis–Januszkiewicz homotopy type*.
An important question arises: to what extent does the isomorphism of the cohomology ring of a space $`X`$ with the face ring $`[K]`$ determine the homotopy type of $`X`$? In other words, for given $`K`$, does there exist a ‘fake’ Davis–Januszkiewicz space, whose cohomology is isomorphic to $`[K]`$, but which is not homotopy equivalent to $`\text{DJ}(K)`$? This question is addressed in . It is shown there \[21, Prop. 5.11\] that if $`[K]`$ is a *complete intersection ring* and $`X`$ is a nilpotent cell complex of finite type whose rational cohomology is isomorphic to $`[K]`$, then $`X`$ is rationally homotopy equivalent to $`\text{DJ}(K)`$. Using the formality of $`\text{DJ}(K)`$, this can be rephrased by saying that the complete intersection face rings are *intrinsically formal* in the sense of Sullivan.
Note that the class of simplicial complexes $`K`$ for which the face ring $`[K]`$ is a complete intersection has a transparent geometrical interpretation: such $`K`$ is a join of simplices and boundaries of simplices.
### 3.3. Coordinate subspace arrangements
Yet another interpretation of the moment-angle complex $`𝒵_K`$ comes from its identification up to homotopy with the complement of the complex coordinate subspace arrangement corresponding to $`K`$. This leads to an application of toric topology in the theory of arrangements, and allows us to describe and effectively calculate the cohomology rings of coordinate subspace arrangement complements and in certain cases identify their homotopy types.
A *coordinate subspace* in $`^m`$ can be written as
(3.4)
$$L_\omega =\{(z_1,\mathrm{},z_m)^m:z_{i_1}=\mathrm{}=z_{i_k}=0\}$$
for some subset $`\omega =\{i_1,\mathrm{},i_k\}[m]`$. Given a simplicial complex $`K`$, we may define the corresponding *coordinate subspace arrangement* $`\{L_\omega :\omega K\}`$ and its *complement*
$$U(K)=^m\underset{\omega K}{}L_\omega .$$
Note that if $`K^{}K`$ is a subcomplex, then $`U(K^{})U(K)`$. It is easy to see \[8, Prop. 8.6\] that the assignment $`KU(K)`$ defines a one-to-one order preserving correspondence between the set of simplicial complexes on $`[m]`$ and the set of coordinate subspace arrangement complements in $`^m`$.
The subset $`U(K)^m`$ is invariant with respect to the coordinatewise $`T^m`$-action. It follows from (3.1) that $`𝒵_KU(K)`$.
###### Proposition 3.8.
There is a $`T^m`$-equivariant deformation retraction
$$U(K)\stackrel{}{}𝒵_K.$$
###### Proof.
In analogy with (3.3), we may write
(3.5)
$$U(K)=\underset{\sigma K}{}U_\sigma ,$$
where
$$U_\sigma :=\{(z_1,\mathrm{},z_m)^m:z_i0\text{ for }i\sigma \}.$$
Then there are obvious homotopy equivalences (deformation retractions)
$$^\sigma \times (0)^{[m]\sigma }U_\sigma \stackrel{}{}B_\sigma (D^2)^\sigma \times (S^1)^{[m]\sigma }.$$
These patch together to get the required map $`U(K)𝒵_K`$. ∎
###### Example 3.9.
1. Let $`K=\mathrm{\Delta }^{m1}`$. Then $`U(K)=^m0`$ (recall that $`𝒵_KS^{2m1}`$ in this case).
2. Let $`K=\{v_1,\mathrm{},v_m\}`$ ($`m`$ points). Then
$$U(K)=^m\underset{1i<jm}{}\{z_i=z_j=0\},$$
the complement to the set of all codimension 2 coordinate planes.
3. More generally, if $`K`$ is the $`i`$-skeleton of $`\mathrm{\Delta }^{m1}`$, then $`U(K)`$ is the complement to the set of all coordinate planes of codimension $`(i+2)`$.
The reader may have noticed a similar pattern in several constructions of toric spaces appeared above; compare (3.1), (3.3) and (3.5). The following general framework was suggested to the author by Neil Strickland in a private communication.
###### Construction 3.10 ($`K`$-power).
Let $`X`$ be a space and $`WX`$ a subspace. For a simplicial complex $`K`$ on $`[m]`$ and $`\sigma K`$, we set
$`(X,W)^\sigma :=\{(x_1,\mathrm{},x_m)X^m:x_jW\text{ for }j\sigma \}`$
and
$`(X,W)^K:={\displaystyle \underset{\sigma K}{}}(X,W)^\sigma ={\displaystyle \underset{\sigma K}{}}({\displaystyle \underset{i\sigma }{}}X\times {\displaystyle \underset{i\sigma }{}}W).`$
We refer to the space $`(X,W)^KX^m`$ as the *$`K`$-power* of $`(X,W)`$. If $`X`$ is a pointed space and $`W=pt`$ is the basepoint, then we shall use the abbreviated notation $`X^K:=(X,pt)^K`$. Examples considered above include $`𝒵_K=(D^2,S^1)^K`$, $`cc(K)=(I^1,S^0)^K`$, $`\text{DJ}(K)=(P^{\mathrm{}})^K`$ and $`U(K)=(,^{})^K`$.
Homotopy theorists would recognise the $`K`$-power as an example of the *colimit* of a diagram of topological spaces over the *face category* of $`K`$ (objects are simplices and morphisms are inclusions). The diagram assigns the space $`(X,W)^\sigma `$ to a simplex $`\sigma `$; its colimit is $`(X,W)^K`$. These observations are further developed and used to construct models of loop spaces of toric spaces as well as for homotopy and homology calculations in and .
### 3.4. Toric varieties, quasitoric manifolds, and torus manifolds
Several important classes of manifolds with torus action emerge as the quotients of moment-angle complexes by appropriate freely acting subtori.
First we give the following characterisation of lsops in the face ring. Let $`K^{n1}`$ be a simplicial complex and $`t_1,\mathrm{},t_n`$ a sequence of degree-two elements in $`𝐤[K]`$. We may write
(3.6)
$$t_i=\lambda _{i1}v_1+\mathrm{}+\lambda _{im}v_m,i=1,\mathrm{},n.$$
For an arbitrary simplex $`\sigma K`$, we have $`K_\sigma =\mathrm{\Delta }^{|\sigma |1}`$ and $`𝐤[K_\sigma ]`$ is the polynomial ring $`𝐤[v_i:i\sigma ]`$ on $`|\sigma |`$ generators. The inclusion $`K_\sigma K`$ induces the *restriction homomorphism* $`r_\sigma `$ from $`𝐤[K]`$ to the polynomial ring, mapping $`v_i`$ identically if $`i\sigma `$ and to zero otherwise.
###### Lemma 3.11.
A degree-two sequence $`t_1,\mathrm{},t_n`$ is an lsop in $`𝐤[K^{n1}]`$ if and only if for every $`\sigma K`$ the elements $`r_\sigma (t_1),\mathrm{},r_\sigma (t_n)`$ generate the positive ideal $`𝐤[v_i:i\sigma ]_+`$.
###### Proof.
Suppose (3.6) is an lsop. For simplicity we denote its image under any restriction homomorphism by the same letters. Then the restriction induces an epimorphism of the quotient rings:
$$𝐤[K]/(t_1,\mathrm{},t_n)𝐤[v_i:i\sigma ]/(t_1,\mathrm{},t_n).$$
Since (3.6) is an lsop, $`𝐤[K]/(t_1,\mathrm{},t_n)`$ is a finitely generated $`𝐤`$-module. Hence, so is $`𝐤[v_i:i\sigma ]/(t_1,\mathrm{},t_n)`$. But the latter can be finitely generated only if $`t_1,\mathrm{},t_n`$ generates $`𝐤[v_i:i\sigma ]_+`$.
The ‘‘if’’ part may be proved by considering the sum of restrictions:
$$𝐤[K]\underset{\sigma K}{}𝐤[v_i:i\sigma ],$$
which turns out to be a monomorphism. See \[6, Th. 5.1.16\] for details. ∎
Obviously, it is enough to consider only restrictions to the maximal simplices in the previous lemma.
Suppose now that $`K`$ is Cohen–Macaulay (e.g. $`K`$ is a sphere triangulation). Then every lsop is a regular sequence (however, for $`𝐤=`$ or a field of finite characteristic an lsop may fail to exist).
Now we restrict to the case $`𝐤=`$ and organise the coefficients in (3.6) into an $`n\times m`$-matrix $`\mathrm{\Lambda }=(\lambda _{ij})`$. For an arbitrary maximal simplex $`\sigma K`$ denote by $`\mathrm{\Lambda }_\sigma `$ the square submatrix formed by the elements $`\lambda _{ij}`$ with $`j\sigma `$. The matrix $`\mathrm{\Lambda }`$ defines a linear map $`^m^n`$ and a homomorphism $`T^mT^n`$. We denote both by $`\lambda `$ and denote the kernel of the latter map by $`T_\mathrm{\Lambda }`$.
###### Theorem 3.12.
The following conditions are equivalent:
* the sequence (3.6) is an lsop in $`[K^{n1}]`$;
* $`det\mathrm{\Lambda }_\sigma =\pm 1`$ for every maximal simplex $`\sigma K`$;
* $`T_\mathrm{\Lambda }T^{mn}`$ and $`T_\mathrm{\Lambda }`$ acts freely on $`𝒵_K`$.
###### Proof.
The equivalence of (a) and (b) is a reformulation of Lemma 3.11. Let us prove the equivalence of (b) and (c). Every isotropy subgroup of the $`T^m`$-action on $`𝒵_K`$ has the form
$$T^\sigma =\{(z_1,\mathrm{},z_m)T^m:z_i=1\text{ if }i\sigma \}$$
for some simplex $`\sigma K`$. Now, (b) is equivalent to the condition $`T_\mathrm{\Lambda }T^\sigma =\{e\}`$ for arbitrary maximal $`\sigma `$, whence the statement follows. ∎
We denote the quotient $`𝒵_K/T_\mathrm{\Lambda }`$ by $`M_K^{2n}(\mathrm{\Lambda })`$, and abbreviate it to $`M_K^{2n}`$ or to $`M^{2n}`$ when the context allows. If $`K`$ is a triangulated sphere, then $`𝒵_K`$ is a manifold, hence, so is $`M_K^{2n}`$. The $`n`$-torus $`T^n=T^m/T_\mathrm{\Lambda }`$ acts on $`M_K^{2n}`$. This construction produces two important classes of $`T^n`$-manifolds as particular examples.
Let $`K=K_P`$ be a polytopal triangulation, dual to the boundary complex of a simple polytope $`P`$. Then the map $`\lambda `$ determined by the matrix $`\mathrm{\Lambda }`$ may be regarded as an assignment of an integer vector to every facet of $`P`$. The map $`\lambda `$ coming from a matrix satisfying the condition of Theorem 3.12(b) was called a *characteristic map* by Davis and Januszkiewicz . We refer to the corresponding quotient $`M_P^{2n}(\mathrm{\Lambda })=𝒵_{K_P}/T_\mathrm{\Lambda }`$ as a *quasitoric manifold* (a toric manifold in the terminology of Davis–Januszkiewicz).
Let us assume further that $`P`$ is realised in $`^n`$ with integer coordinates of vertices, so we can write
(3.7)
$$P^n=\{\text{x}^n:\text{l}_i,\text{x}a_i,i=1,\mathrm{},m\},$$
where $`\text{l}_i`$ are inward pointing normals to the facets of $`P^n`$ (we may further assume these vectors to be primitive), and $`a_i`$. Let $`\mathrm{\Lambda }`$ be the matrix formed by the column vectors $`\text{l}_i`$, $`i=1,\mathrm{},m`$. Then $`𝒵_{K_P}/T_\mathrm{\Lambda }`$ can be identified with the *projective toric variety* determined by the polytope $`P`$. The condition of Theorem 3.12(b) is equivalent to the requirement that the toric variety is non-singular. Thereby a non-singular projective toric variety is a quasitoric manifold (but there are many quasitoric manifolds which are not toric varieties).
We also note that smooth projective toric varieties provide examples of *symplectic* $`2n`$-dimensional manifolds with *Hamiltonian* $`T^n`$-action. These symplectic manifolds can be obtained via the process of *symplectic reduction* from the standard Hamiltonian $`T^m`$-action on $`^m`$. A choice of an $`(mn)`$-dimensional toric subgroup provides a *moment map* $`\mu :^m^{mn}`$, and the corresponding moment-angle complex $`𝒵_{K_P}`$ can be identified with the level surface $`\mu ^1(a)`$ of the moment map for any of its regular values $`a`$. The details of this construction can be found in \[8, p. 130\].
Finally, we mention that if $`K`$ is an arbitrary (not necessarily polytopal) triangulation of sphere, then the manifold $`M_K^{2n}(\mathrm{\Lambda })`$ is a *torus manifold* in the sense of Hattori–Masuda . The corresponding multi-fan has $`K`$ as the underlying simplicial complex. This particular class of torus manifolds has many interesting properties.
## 4. Cohomology of moment-angle complexes
The main result of this section (Theorem 4.7) identifies the *integral* cohomology algebra of the moment-angle complex $`𝒵_K`$ with the $`\mathrm{Tor}`$-algebra of the face ring of the simplicial complex $`K`$. Over the rationals this result was proved in by studying the Eilenberg–Moore spectral sequence of the fibration $`𝒵_K\text{DJ}(K)BT^m`$; a more detailed account of applications of the Eilenberg–Moore spectral sequence to toric topology can be found in . The new proof, which works with integer coefficients as well, relies upon a construction of a special cellular decomposition of $`𝒵_K`$ and subsequent analysis of the corresponding cellular cochains.
One of the key ingredients here is a specific cellular approximation of the diagonal map $`\mathrm{\Delta }:𝒵_K𝒵_K\times 𝒵_K`$. Cellular cochains do not admit a functorial associative multiplication because a proper cellular diagonal approximation does not exist in general. The construction of moment-angle complexes is given by a functor from the category of simplicial complexes to the category of spaces with a torus action. We show that in this special case the cellular approximation of the diagonal is functorial with respect to those maps of moment-angle complexes which are induced by simplicial maps. The corresponding cellular cochain algebra is isomorphic to a quotient of the Koszul complex for $`𝐤[K]`$ by an acyclic ideal, and its cohomology is isomorphic to the $`\mathrm{Tor}`$-algebra. The proofs have been sketched in ; here we follow the more detailed exposition of . Another proof of Theorem 4.7 follows from a recent independent work of M. Franz \[12, Th. 1.2\].
### 4.1. Cell decomposition
The polydisc $`(D^2)^m`$ has a cell decomposition in which each $`D^2`$ is subdivided into cells 1, $`T`$ and $`D`$ of dimensions $`0`$, $`1`$ and $`2`$ respectively, see Figure 3.
Each cell of this complex is a product of cells of 3 different types and we encode it by a word $`𝒯\{D,T,1\}^m`$ in a three-letter alphabet. Assign to each pair of subsets $`\sigma ,\omega [m]`$, $`\sigma \omega =\mathrm{}`$, the word $`𝒯(\sigma ,\omega )`$ which has the letter $`D`$ on the positions indexed by $`\sigma `$ and letter $`T`$ on the positions with indices from $`\omega `$.
###### Lemma 4.1.
$`𝒵_K`$ is a cellular subcomplex of $`(D^2)^m`$. A cell $`𝒯(\sigma ,\omega )(D^2)^m`$ belongs to $`𝒵_K`$ if and only if $`\sigma K`$.
###### Proof.
We have $`𝒵_K=_{\sigma K}B_\sigma `$ and each $`B_\sigma `$ is the closure of the cell $`𝒯(\sigma ,[m]\sigma )`$. ∎
Therefore, we can consider the cellular cochain complex $`C^{}(𝒵_K)`$, which has an additive basis consisting of the cochains $`𝒯(\sigma ,\omega )^{}`$. It has a natural bigrading defined by
$$bideg𝒯(\sigma ,\omega )^{}=(|\omega |,2|\sigma |+2|\omega |),$$
so $`bidegD=(0,2)`$, $`bidegT=(1,2)`$ and $`bideg1=(0,0)`$. Moreover, since the cellular differential does not change the second grading, $`C^{}(𝒵_K)`$ splits into the sum of its components having fixed second degree:
$$C^{}(𝒵_K)=\underset{j=1}{\overset{m}{}}C^{,2j}(𝒵_K).$$
The cohomology of $`𝒵_K`$ thereby acquires an additional grading, and we may define the *bigraded Betti numbers* $`b^{i,2j}(𝒵_K)`$ by
$$b^{i,2j}(𝒵_K):=rankH^{i,2j}(𝒵_K),i,j=1,\mathrm{},m.$$
For the ordinary Betti numbers we have $`b^k(𝒵_K)=_{2ji=k}b^{i,2j}(𝒵_K)`$.
###### Lemma 4.2.
Let $`\phi :K_1K_2`$ be a simplicial map between simplicial complexes on the sets $`[m_1]`$ and $`[m_2]`$ respectively. Then there is an equivariant cellular map $`\phi _𝒵:𝒵_{K_1}𝒵_{K_2}`$ covering the induced map $`|coneK_1^{}||coneK_2^{}|`$.
###### Proof.
Define a map of polydisks
$$\phi _D:(D^2)^{m_1}(D^2)^{m_2},(z_1,\mathrm{},z_{m_1})(w_1,\mathrm{},w_{m_2}),$$
where
$$w_j:=\underset{i\phi ^1(j)}{}z_i,j=1,\mathrm{},m_2$$
(we set $`w_j=1`$ if $`\phi ^1(j)=\mathrm{}`$). Assume $`\tau K_1`$. In the notation of (3.1), we have $`\phi _D(B_\tau )B_{\phi (\tau )}`$. Since $`\phi `$ is a simplicial map, $`\phi (\tau )K_2`$ and $`B_{\phi (\tau )}𝒵_{K_2}`$, so the restriction of $`\phi _D`$ to $`𝒵_{K_1}`$ is the required map. ∎
###### Corollary 4.3.
The correspondence $`K𝒵_K`$ gives rise to a functor from the category of simplicial complexes and simplicial maps to the category of spaces with torus actions and equivariant maps. It induces a natural transformation between the simplicial cochain functor of $`K`$ and the cellular cochain functor of $`𝒵_K`$.
We also note that the maps respect the bigrading, so the bigraded Betti numbers are also functorial.
### 4.2. Koszul algebras
Our algebraic model for the cellular cochains of $`𝒵_K`$ is obtained by taking the quotient of the Koszul algebra $`[\mathrm{\Lambda }[u_1,\mathrm{},u_m]𝐤[K],d]`$ from Lemma 2.11 by a certain acyclic ideal. Namely, we introduce a factor algebra
$$R^{}(K):=\mathrm{\Lambda }[u_1,\mathrm{},u_m][K]/(v_i^2=u_iv_i=0,i=1,\mathrm{},m),$$
where the differential and bigrading are as in (2.4). Let
$$\varrho :\mathrm{\Lambda }[u_1,\mathrm{},u_m][K]R^{}(K)$$
be the quotient projection. The algebra $`R^{}(K)`$ has a finite additive basis consisting of the monomials of the form $`u_\omega v_\sigma `$ where $`\omega [m]`$, $`\sigma K`$ and $`\omega \sigma =\mathrm{}`$ (remember that we are using the notation $`u_\omega =u_{i_1}\mathrm{}u_{i_k}`$ for $`\omega =\{i_1,\mathrm{},i_k\}`$). Therefore, we have an additive inclusion (a monomorphism of bigraded differential modules)
$$\iota :R^{}(K)\mathrm{\Lambda }[u_1,\mathrm{},u_m][K]$$
which satisfies $`\varrho \iota =\mathrm{id}`$.
The following statement shows that the finite-dimensional quotient $`R^{}(K)`$ has the same cohomology as the Koszul algebra.
###### Lemma 4.4.
The quotient map $`\varrho :\mathrm{\Lambda }[u_1,\mathrm{},u_m][K]R^{}(K)`$ induces an isomorphism in cohomology.
###### Proof.
The argument is similar to that used in the proof of the acyclicity of the Koszul resolution. We construct a cochain homotopy between the maps $`\mathrm{id}`$ and $`\iota \varrho `$ from $`\mathrm{\Lambda }[u_1,\mathrm{},u_m][K]`$ to itself, that is, a map $`s`$ satisfying
(4.1)
$$ds+sd=\mathrm{id}\iota \varrho .$$
First assume that $`K=\mathrm{\Delta }^{m1}`$. We denote the corresponding bigraded algebra $`\mathrm{\Lambda }[u_1,\mathrm{},u_m][K]`$ by
(4.2)
$$E=E_m:=\mathrm{\Lambda }[u_1,\mathrm{},u_m][v_1,\mathrm{},v_m],$$
while $`R^{}(K)`$ is isomorphic to
(4.3)
$$\left(\mathrm{\Lambda }[u][v]/(v^2=uv=0)\right)^m=R^{}(\mathrm{\Delta }^0)^m.$$
For $`m=1`$, the map $`s_1:E^{0,}=𝐤[v]E^{1,}`$ given by
$$s_1(a_0+a_1v+\mathrm{}+a_jv^j)=(a_2v+a_3v^2+\mathrm{}+a_jv^{j1})u$$
is a cochain homotopy. Indeed, we can write an element of $`E`$ as either $`x`$ or $`xu`$ with $`x=a_0+a_1v+\mathrm{}+a_jv^jE^{0,2j}`$. In the former case, $`ds_1x=xa_0a_1v=x\iota \varrho x`$ and $`s_1dx=0`$. In the latter case, $`xuE^{1,2j}`$, then $`ds_1(xu)=0`$ and $`s_1d(xu)=xua_0u=xu\iota \varrho (xu)`$. In both cases (4.1) holds. Now we may assume by induction that for $`m=k1`$ there is a cochain homotopy operator $`s_{k1}:E_{k1}E_{k1}`$. Since $`E_k=E_{k1}E_1`$, $`\varrho _k=\varrho _{k1}\varrho _1`$ and $`\iota _k=\iota _{k1}\iota _1`$, a direct check shows that the map
$$s_k=s_{k1}\mathrm{id}+\iota _{k1}\varrho _{k1}s_1$$
is a cochain homotopy between $`\mathrm{id}`$ and $`\iota _k\varrho _k`$, which finishes the proof for $`K=\mathrm{\Delta }^{m1}`$.
In the case of arbitrary $`K`$ the algebras $`\mathrm{\Lambda }[u_1,\mathrm{},u_m][K]`$ are $`R^{}(K)`$ are obtained from (4.2) and (4.3) respectively by factoring out the Stanley–Reisner ideal $`_K`$. This factorisation does not affect the properties of the constructed map $`s`$, which finishes the proof. ∎
Now comparing the additive structure of $`R^{}(K)`$ with that of the cellular cochains $`C^{}(K)`$, we see that the two coincide:
###### Lemma 4.5.
The map
$`g:R^{}(K)`$ $`C^{}(𝒵_K),`$
$`u_\omega v_\sigma `$ $`𝒯(\sigma ,\omega )^{}`$
is an isomorphism of bigraded differential modules. In particular, we have an additive isomorphism
$$H[R^{}(K)]H^{}(𝒵_K).$$
Having identified the algebra $`R^{}`$ with the cellular cochains of $`𝒵_K`$, we can also interpret the cohomology isomorphism from Lemma 4.4 topologically. To do this we shall identify the Koszul algebra $`\mathrm{\Lambda }[u_1,\mathrm{},u_m][K]`$ with the cellular cochains of a space homotopy equivalent to $`𝒵_K`$.
Let $`S^{\mathrm{}}`$ be an infinite-dimensional sphere obtained as a direct limit (union) of standardly embedded odd-dimensional spheres. The space $`S^{\mathrm{}}`$ is contractible and has a cell decomposition with one cell in every dimension. The boundary of an even-dimensional cell is the closure of the appropriate odd-dimensional cell, while the boundary of an odd cell is zero. The 2-skeleton of this cell decomposition is a 2-disc decomposed as shown on Figure 3, while the 1-skeleton is the circle $`S^1S^{\mathrm{}}`$. The cellular cochain complex of $`S^{\mathrm{}}`$ can be identified with the algebra
$$\mathrm{\Lambda }[u][v],\mathrm{deg}u=1,\mathrm{deg}v=2,du=v,dv=0.$$
From the obvious functorial properties of Construction 3.10 we obtain a deformation retraction
$$𝒵_K=(D^2,S^1)^K(S^{\mathrm{}},S^1)^K(D^2,S^1)^K$$
onto a cellular subcomplex.
The cellular cochains of the $`K`$-power $`(S^{\mathrm{}},S^1)^K`$ can be identified with the Koszul algebra $`\mathrm{\Lambda }[u_1,\mathrm{},u_m][K]`$. Since $`𝒵_K(S^{\mathrm{}},S^1)^K`$ is a deformation retract, the cellular cochain map
$$\mathrm{\Lambda }[u_1,\mathrm{},u_m][K]=C^{}\left((S^{\mathrm{}},S^1)^K\right)C^{}(𝒵_K)=R^{}(K),$$
induces an isomorphism in cohomology. In fact, the algebraic homotopy map $`s`$ constructed in the proof of Lemma 4.4 is the map induced on the cochains by the topological homotopy.
### 4.3. Cellular cochain algebras
Here we introduce a multiplication for cellular cochains of $`𝒵_K`$ and establish a ring isomorphism in Lemma 4.5. This task runs into a complication because cellular cochains in general do not carry a functorial associative multiplication; the classical definition of the cohomology multiplication involves a diagonal map, which is not cellular. However, in our case there is a way to construct a canonical cellular approximation of the diagonal map $`\mathrm{\Delta }:𝒵_K𝒵_K\times 𝒵_K`$ in such a way that the resulting multiplication in cellular cochains coincides with that in $`R^{}(K)`$.
The standard definition of the multiplication in cohomology of a cell complex $`X`$ via cellular cochains is as follows. Consider a composite map of cellular cochain complexes:
(4.4)
$$\begin{array}{ccccc}C^{}(X)C^{}(X)& \stackrel{\times }{}& C^{}(X\times X)& \stackrel{\stackrel{~}{\mathrm{\Delta }}^{}}{}& C^{}(X).\end{array}$$
Here the map $`\times `$ assigns to a cellular cochain $`c_1c_2C^{q_1}(X)C^{q_2}(X)`$ the cochain $`c_1\times c_2C^{q_1+q_2}(X\times X)`$ whose value on a cell $`e_1\times e_2X\times X`$ is $`(1)^{q_1q_2}c_1(e_1)c_2(e_2)`$. The map $`\stackrel{~}{\mathrm{\Delta }}^{}`$ is induced by a cellular approximation $`\stackrel{~}{\mathrm{\Delta }}`$ of the diagonal map $`\mathrm{\Delta }:XX\times X`$. In cohomology, the map (4.4) induces a multiplication $`H^{}(X)H^{}(X)H^{}(X)`$ which does not depend on a choice of cellular approximation and is functorial. However, the map (4.4) is not itself functorial because of the arbitrariness in the choice of a cellular approximation.
In the special case $`X=𝒵_K`$ we may apply the following construction. Consider a map $`\stackrel{~}{\mathrm{\Delta }}:D^2D^2\times D^2`$, defined in polar coordinates $`z=\rho e^{i\phi }D^2`$, $`0\rho 1`$, $`0\phi <2\pi `$ as follows:
$$\rho e^{i\phi }\{\begin{array}{cc}(1+\rho (e^{2i\phi }1),1)\hfill & \text{ for }0\phi \pi ,\hfill \\ (1,1+\rho (e^{2i\phi }1))\hfill & \text{ for }\pi \phi <2\pi .\hfill \end{array}$$
This is a cellular map taking $`D^2`$ to $`D^2\times D^2`$ and homotopic to the diagonal $`\mathrm{\Delta }:D^2D^2\times D^2`$ in the class of such maps. Taking an $`m`$-fold product, we obtain a cellular approximation
$$\stackrel{~}{\mathrm{\Delta }}:(D^2)^m(D^2)^m\times (D^2)^m$$
which restricts to a cellular approximation for the diagonal map of $`𝒵_K`$ for arbitrary $`K`$, as described by the following commutative diagram:
$$\begin{array}{ccc}𝒵_K& & (D^2)^m\\ \stackrel{~}{\mathrm{\Delta }}& & \stackrel{~}{\mathrm{\Delta }}& & \\ 𝒵_K\times 𝒵_K& & (D^2)^m\times (D^2)^m\end{array}.$$
Note that this diagonal approximation is functorial with respect to those maps $`𝒵_{K_1}𝒵_{K_2}`$ of moment-angle complexes that are induced by simplicial maps $`K_1K_2`$ (see Lemma 4.2).
###### Lemma 4.6.
The cellular cochain algebra $`C^{}(𝒵_K)`$ defined by the diagonal approximation $`\stackrel{~}{\mathrm{\Delta }}:𝒵_K𝒵_K\times 𝒵_K`$ and (4.4) is isomorphic to $`R^{}(K)`$. Therefore, we get an isomorphism of the cohomology algebras:
$$H[R^{}(K)]H^{}(𝒵_K;).$$
###### Proof.
We first consider the case $`K=\mathrm{\Delta }^0`$, that is, $`𝒵_K=D^2`$. The cellular cochain complex of $`D^2`$ is additively generated by the cochains $`1C^0(D^2)`$, $`T^{}C^1(D^2)`$ and $`D^{}C^2(D^2)`$ dual to the corresponding cells, see Figure 3. The multiplication defined in $`C^{}(D^2)`$ by (4.4) is trivial, so we get a multiplicative isomorphism
$$R^{}(\mathrm{\Delta }^0)=\mathrm{\Lambda }[u][v]/(v^2=uv=0)C^{}(D^2).$$
Now, for $`K=\mathrm{\Delta }^{m1}`$ we obtain a multiplicative isomorphism
$$f:R^{}(\mathrm{\Delta }^{m1})=\mathrm{\Lambda }[u_1,\mathrm{},u_m][v_1,\mathrm{},v_m]/(v_i^2=u_iv_i=0)C^{}((D^2)^m)$$
by taking the tensor product. Since $`𝒵_K(D^2)^m`$ is a cell subcomplex for arbitrary $`K`$ we obtain a multiplicative map $`q:C^{}((D^2)^m)C^{}(𝒵_K)`$. Now consider the commutative diagram
$$\begin{array}{ccc}R^{}(\mathrm{\Delta }^{m1})& \stackrel{f}{}& C^{}((D^2)^m)\\ p& & q& & \\ R^{}(K)& \stackrel{g}{}& C^{}(𝒵_K).\end{array}$$
Here the maps $`p`$, $`f`$ and $`q`$ are multiplicative, while $`g`$ is an additive isomorphism by Lemma 4.5. Take $`\alpha ,\beta R^{}(K)`$. Since $`p`$ is onto, we have $`\alpha =p(\alpha ^{})`$ and $`\beta =p(\beta ^{})`$. Then
$$g(\alpha \beta )=gp(\alpha ^{}\beta ^{})=qf(\alpha ^{}\beta ^{})=gp(\alpha ^{})gp(\beta ^{})=g(\alpha )g(\beta ),$$
and $`g`$ is also a multiplicative isomorphism, which finishes the proof. ∎
Combining the results of Lemmas 2.11, 2.12, 4.4 and 4.6, we come to the main result of this section.
###### Theorem 4.7.
There is an isomorphism, functorial in $`K`$, of bigraded algebras
$$H^,(𝒵_K;)\mathrm{Tor}_{[v_1,\mathrm{},v_m]}([K],)H[\mathrm{\Lambda }[u_1,\mathrm{},u_m][K],d],$$
where the bigrading and the differential in the last algebra are defined by (2.4).
As an illustration, we give two examples of particular cohomology calculations, which have a transparent geometrical interpretation. More examples of calculations may be found in .
###### Example 4.8.
1. Let $`K=\mathrm{\Delta }^{m1}`$. Then
$$[K]=[v_1,\mathrm{},v_m]/(v_1\mathrm{}v_m).$$
The fundamental class of $`𝒵_KS^{2m1}`$ is represented by the bideg $`(1,2m)`$ cocycle $`u_1v_2v_3\mathrm{}v_m\mathrm{\Lambda }[u_1,\mathrm{},u_m][K]`$.
2. Let $`K=\{v_1,\mathrm{},v_m\}`$ ($`m`$ points). Then $`𝒵_K`$ is homotopy equivalent to the complement in $`^m`$ to the set of all codimension-two coordinate planes, see Example 3.9. Then
$$[K]=[v_1,\mathrm{},v_m]/(v_iv_j,ij).$$
The subspace of cocycles in $`R^{}(K)`$ is generated by
$$v_{i_1}u_{i_2}u_{i_3}\mathrm{}u_{i_k},k2\text{ and }i_pi_q\text{ for }pq,$$
and has dimension $`m\left(\genfrac{}{}{0pt}{}{m1}{k1}\right)`$. The subspace of coboundaries is generated by the elements of the form
$$d(u_{i_1}\mathrm{}u_{i_k})$$
and is $`\left(\genfrac{}{}{0pt}{}{m}{k}\right)`$-dimensional. Therefore
$$\begin{array}{c}dimH^0(𝒵_K)=1,\hfill \\ dimH^1(𝒵_K)=H^2(U(K))=0,\hfill \\ dimH^{k+1}(𝒵_K)=m\left(\genfrac{}{}{0pt}{}{m1}{k1}\right)\left(\genfrac{}{}{0pt}{}{m}{k}\right)=(k1)\left(\genfrac{}{}{0pt}{}{m}{k}\right),2km,\hfill \end{array}$$
and multiplication in the cohomology of $`𝒵_K`$ is trivial. Note that in general multiplication in the cohomology of $`𝒵_K`$ is far from being trivial; for example if $`K`$ is a sphere triangulation then $`𝒵_K`$ is a manifold by Lemma 3.3.
The above cohomology calculation suggests that the complement of the subspace arrangement from the previous example is homotopy equivalent to a wedge of spheres. This is indeed the case, as the following theorem shows.
###### Theorem 4.9 (Grbić–Theriault ).
The complement of the set of all codimension-two coordinate subspaces in $`^m`$ has the homotopy type of the wedge of spheres
$$\underset{k=2}{\overset{m}{}}(k1)\left(\genfrac{}{}{0pt}{}{m}{k}\right)S^{k+1}.$$
The proof is based on an analysis of the homotopy fibre of the inclusion $`\text{DJ}(K)BT^m`$, which is homotopy equivalent to $`𝒵_K`$ (or $`U(K)`$) by Corollary 3.7. We shall return to coordinate subspace arrangements once again in the next section.
## 5. Applications to combinatorial commutative algebra
### 5.1. A multiplicative version of Hochster’s theorem
As a first application we give a proof of a generalisation of Hochster’s theorem (Theorem 2.9) obtained by Baskakov in .
The bigraded structure in the cellular cochains of $`𝒵_K`$ can be further refined as
$$C^{}(𝒵_K)=\underset{\omega [m]}{}C^{,\mathrm{\hspace{0.17em}2}\omega }(𝒵_K)$$
where $`C^{,\mathrm{\hspace{0.17em}2}\omega }(𝒵_K)`$ is the subcomplex generated by the cochains $`𝒯(\sigma ,\omega \sigma )^{}`$ with $`\sigma \omega `$ and $`\sigma K`$. Thus, $`C^{}(𝒵_K)`$ now becomes a $`^m`$-graded module, and the bigraded cohomology groups decompose accordingly as
(5.1)
$$H^{i,\mathrm{\hspace{0.17em}2}j}(𝒵_K)=\underset{\omega [m]:|\omega |=j}{}H^{i,\mathrm{\hspace{0.17em}2}\omega }(𝒵_K)$$
where $`H^{i,\mathrm{\hspace{0.17em}2}\omega }(𝒵_K):=H^i[C^{,\mathrm{\hspace{0.17em}2}\omega }(𝒵_K)]`$.
Given two simplicial complexes $`K_1`$ and $`K_2`$ with vertex sets $`V_1`$ and $`V_2`$ respectively, their *join* is the following complex on $`V_1V_2`$:
$$K_1K_2:=\{\sigma V_1V_2:\sigma =\sigma _1\sigma _2,\sigma _1K_1,\sigma _2K_2\}.$$
Now we introduce a multiplication in the sum
$$\underset{\omega [m]}{\underset{p1,}{}}\stackrel{~}{H}^p(K_\omega )$$
where $`K_\omega `$ is the full subcomplex and $`\stackrel{~}{H}^1(\mathrm{})=`$, as follows. Take two elements $`\alpha \stackrel{~}{H}^p(K_{\omega _1})`$ and $`\beta \stackrel{~}{H}^q(K_{\omega _2})`$. Assume that $`\omega _1\omega _2=\mathrm{}`$. Then we have an inclusion of subcomplexes
$$i:K_{\omega _1\omega _2}=K_{\omega _1}K_{\omega _2}K_{\omega _1}K_{\omega _2}$$
and an isomorphism of reduced simplicial cochains
$$f:\stackrel{~}{C}^p(K_{\omega _1})\stackrel{~}{C}^q(K_{\omega _2})\stackrel{}{}\stackrel{~}{C}^{p+q+1}(K_{\omega _1}K_{\omega _2}).$$
Now set
$$\alpha \beta :=\{\begin{array}{cc}0,\hfill & \omega _1\omega _2\mathrm{},\hfill \\ i^{}f(ab)\stackrel{~}{H}^{p+q+1}(K_{\omega _1\omega _2}),\hfill & \omega _1\omega _2=\mathrm{}.\hfill \end{array}$$
###### Theorem 5.1 (Baskakov \[3, Th. 1\]).
There are isomorphisms
$$\stackrel{~}{H}^p(K_\omega )\stackrel{}{}H^{p+1|\omega |,2\omega }(𝒵_K)$$
which are functorial with respect to simplicial maps and induce a ring isomorphism
$$\gamma :\underset{\omega [m]}{\underset{p1,}{}}\stackrel{~}{H}^p(K_\omega )\stackrel{}{}H^{}(𝒵_K).$$
###### Proof.
Define a map of cochain complexes
$$\stackrel{~}{C}^{}(K_\omega )C^{+1|\omega |,2\omega }(𝒵_K),\sigma ^{}𝒯(\sigma ,\omega \sigma )^{}.$$
It is a functorial isomorphism by observation, whence the isomorphism of the cohomology groups follows.
The statement about the ring isomorphism follows from the isomorphism $`H^{}(𝒵_K)H[R^{}(K)]`$ established in Lemma 4.5 and analysing the ring structure in $`R^{}(K)`$. ∎
###### Corollary 5.2.
There is an isomorphism
$$H^{i,2j}(𝒵_K)\underset{\omega [m]:|\omega |=j}{}\stackrel{~}{H}^{ji1}(K_\omega ).$$
As a further corollary we obtain Hochster’s theorem (Theorem 2.9):
$$\mathrm{Tor}_{[v_1,\mathrm{},v_m]}^{i,}([K],)\underset{\omega [m]}{}\stackrel{~}{H}^{|\omega |i1}(K_\omega ).$$
### 5.2. Alexander duality and coordinate subspace arrangements revisited
The multiplicative version of Hochster’s can also be applied to cohomology calculations of subspace arrangement complements.
A coordinate subspace can be defined either by setting some coordinates to zero as in (3.4), or as the linear span of a subset of the standard basis in $`^m`$. This gives an alternative way to parametrise coordinate subspace arrangements by simplicial complexes. Namely, we can write
$$\{L_\omega :\omega K\}=\{spane_{i_1},\mathrm{},e_{i_k}:\{i_1,\mathrm{},i_k\}\widehat{K}\}$$
where $`\widehat{K}`$ is the simplicial complex given by
$$\widehat{K}:=\{\omega [m]:[m]\omega K\}.$$
It is called the *dual complex* of $`K`$. The cohomology of full subcomplexes in $`K`$ is related to the homology of links in $`\widehat{K}`$ by means of the following combinatorial version of the Alexander duality theorem.
###### Theorem 5.3 (Alexander duality).
Let $`K\mathrm{\Delta }^{m1}`$ be a simplicial complex on the set $`[m]`$ and $`\sigma K`$, that is, $`\widehat{\sigma }=[m]\sigma \widehat{K}`$. Then there are isomorphisms
$$\stackrel{~}{H}_j(K_\sigma )\stackrel{~}{H}^{|\sigma |3j}(link_{\widehat{K}}\widehat{\sigma }).$$
In particular, for $`\sigma =[m]`$ we get
$$\stackrel{~}{H}_j(K)\stackrel{~}{H}^{m3j}(\widehat{K}),1jm2.$$
A proof can be found in \[9, §2.2\]. Using the duality between the full subcomplexes of $`K`$ and links of $`\widehat{K}`$ we can reformulate the cohomology calculation of $`U(K)`$ as follows.
###### Proposition 5.4.
We have
$$\stackrel{~}{H}_i(U(K))\underset{\sigma \widehat{K}}{}\stackrel{~}{H}^{2m2|\sigma |i2}(link_{\widehat{K}}\sigma ).$$
###### Proof.
From Proposition 3.8 and Corollary 5.2 we obtain
$$H_p(U(K))=\underset{\tau [m]}{}\stackrel{~}{H}_{p|\tau |1}(K_\tau ).$$
Nonempty simplices $`\tau K`$ do not contribute to the above sum, since the corresponding subcomplexes $`K_\tau `$ are contractible. Since $`\stackrel{~}{H}^1(\mathrm{})=𝐤`$ the empty subset of $`[m]`$ only contributes $`𝐤`$ to $`H^0(U(K))`$. Hence we may rewrite the above formula as
$$\stackrel{~}{H}_p(U(K))=\underset{\tau K}{}\stackrel{~}{H}_{p|\tau |1}(K_\tau ).$$
Using Theorem 5.3, we calculate
$$\stackrel{~}{H}_{p|\tau |1}(K_\tau )=\stackrel{~}{H}^{|\tau |3p+|\tau |+1}(link_{\widehat{K}}\widehat{\tau })=\stackrel{~}{H}^{2m2|\widehat{\tau }|p2}(link_{\widehat{K}}\widehat{\tau }),$$
where $`\widehat{\tau }=[m]\tau `$ is a simplex in $`\widehat{K}`$, as required. ∎
Proposition 5.4 is a particular case of the well-known *Goresky–Macpherson formula* \[15, Part III\], which calculates the dimensions of the (co)homology groups of an arbitrary subspace arrangement in terms of its *intersection poset* (which coincides with the poset of faces of $`\widehat{K}`$ in the case of coordinate arrangements). We see that the study of moment-angle complexes not only allows us to retrieve the multiplicative structure of the cohomology of complex coordinate subspace arrangement complements, but also connects two seemingly unrelated results, the Goresky–Macpherson formula from the theory of arrangements and Hochester’s formula from combinatorial commutative algebra.
### 5.3. Massey products in the cohomology of $`𝒵_K`$
Here we address the question of existence of non-trivial Massey products in the Koszul complex
$$[\mathrm{\Lambda }[u_1,\mathrm{},u_m][K],d]$$
of the face ring. Massey products constitute a series of higher-order operations (or *brackets*) in the cohomology of a differential graded algebra, with the second-order operation coinciding with the cohomology multiplication, while the higher-order brackets are only defined for certain tuples of cohomology classes. A geometrical approach to constructing nontrivial triple Massey products in the Koszul complex of the face ring has been developed by Baskakov in as an extension of the cohomology calculation in Theorem 5.1. It is well-known that non-trivial higher Massey products obstruct the *formality* of a differential graded algebra, which in our case leads to a family on nonformal moment-angle manifolds $`𝒵_K`$.
Massey products in the cohomology of the Koszul complex of a local ring $`R`$ were studied by Golod in connection with the calculation of the Poincaré series of $`\mathrm{Tor}_R(𝐤,𝐤)`$. The main result of Golod is a calculation of the Poincaré series for the class of rings with vanishing Massey products in the Koszul complex (including the cohomology multiplication). Such rings were called *Golod* in , where the reader can find a detailed exposition of Golod’s theorem together with several further applications. The Golod property of face rings was studied in , where several combinatorial criteria for Golodness were given.
The difference between our situation and that of Golod is that we are mainly interested in the cohomology of the Koszul complex for the face ring of a sphere triangulation $`K`$. The corresponding face ring $`𝐤[K]`$ does not qualify for Golodness, as the corresponding moment-angle complex $`𝒵_K`$ is a manifold, and therefore, the cohomology of the Koszul complex of $`𝐤[K]`$ must possess many non-trivial products. Our approach aims at identifying a class of simplicial complexes with non-trivial cohomology product but vanishing higher-order Massey operations in the cohomology of the Koszul complex.
Let $`K_i`$ be a triangulation of a sphere $`S^{n_i1}`$ with $`|V_i|=m_i`$ vertices, $`i=1,2,3`$. Set $`m:=m_1+m_2+m_3`$, $`n:=n_1+n_2+n_3`$, and
$$\begin{array}{cc}K:=K_1K_2K_3,𝒵_K=𝒵_{K_1}\times 𝒵_{K_2}\times 𝒵_{K_3}.\hfill & \end{array}$$
Note that $`K`$ is a triangulation of $`S^{n1}`$ and $`𝒵_K`$ is an $`(m+n)`$-manifold.
Given $`\sigma K`$, the *stellar subdivision* of $`K`$ at $`\sigma `$ is obtained by replacing the star of $`\sigma `$ by the cone over its boundary:
$$\zeta _\sigma (K)=(Kstar_K\sigma )(conestar_K\sigma ).$$
Now choose maximal simplices $`\sigma _1K_1`$, $`\sigma _2^{},\sigma _2^{\prime \prime }K_2`$ such that $`\sigma _2^{}\sigma _2^{\prime \prime }=\mathrm{}`$, and $`\sigma _3K_3`$. Set
$$\stackrel{~}{K}:=\zeta _{\sigma _1\sigma _2^{}}(\zeta _{\sigma _2^{\prime \prime }\sigma _3}(K)).$$
Then $`\stackrel{~}{K}`$ is a triangulation of $`S^{n1}`$ with $`m+2`$ vertices. Take generators
$$\beta _i\stackrel{~}{H}^{n_i1}(\stackrel{~}{K}_{V_i})\stackrel{~}{H}^{n_i1}(S^{n_i1}),i=1,2,3,$$
where $`\stackrel{~}{K}_{V_i}`$ is the restriction of $`\stackrel{~}{K}`$ to the vertex set of $`K_i`$, and set
$$\alpha _i:=\gamma (\beta _i)H^{n_im_i,2m_i}(𝒵_{\stackrel{~}{K}})H^{m_i+n_i}(𝒵_{\stackrel{~}{K}}),$$
where $`\gamma `$ is the isomorphism from Theorem 5.1. Then
$$\beta _1\beta _2\stackrel{~}{H}^{n_1+n_21}(\stackrel{~}{K}_{V_1V_2})\stackrel{~}{H}^{n_1+n_21}(S^{n_1+n_21}\text{pt})=0,$$
and therefore, $`\alpha _1\alpha _2=\gamma (\beta _1\beta _2)=0`$, and similarly $`\alpha _2\alpha _3=0`$. In these circumstances the triple Massey product $`\alpha _1,\alpha _2,\alpha _3H^{m+n1}(𝒵_{\stackrel{~}{K}})`$ is defined. Recall that $`\alpha _1,\alpha _2,\alpha _3`$ is the set of cohomology classes represented by the cocycles $`(1)^{\mathrm{deg}a_1+1}a_1f+ea_3`$ where $`a_i`$ is a cocycle representing $`\alpha _i`$, $`i=1,2,3`$, while $`e`$ and $`f`$ are cochains satisfying $`de=a_1a_2`$, $`df=a_2a_3`$. A Massey product is called *trivial* if it contains zero.
###### Theorem 5.5.
The triple Massey product
$$\alpha _1,\alpha _2,\alpha _3H^{m+n1}(𝒵_{\stackrel{~}{K}})$$
in the cohomology of $`(m+n+2)`$-manifold $`𝒵_{\stackrel{~}{K}}`$ is non-trivial.
###### Proof.
Consider the subcomplex of $`\stackrel{~}{K}`$ consisting of those two new vertices added to $`K`$ in the process of stellar subdivision. By Lemma 4.2, the inclusion of this subcomplex induces an embedding of a 3-dimensional sphere $`S^3𝒵_{\stackrel{~}{K}}`$. Since the two new vertices are not joined by an edge in $`𝒵_{\stackrel{~}{K}}`$, the embedded 3-sphere defines a non-trivial class in $`H^3(𝒵_{\stackrel{~}{K}})`$. By construction the dual cohomology class is contained in the Massey product $`\alpha _1,\alpha _2,\alpha _3`$. On the other hand, this Massey product is defined up to elements from the subspace
$$\alpha _1H^{m_2+m_3+n_2+n_31}(𝒵_{\stackrel{~}{K}})+\alpha _3H^{m_1+m_2+n_1+n_21}(𝒵_{\stackrel{~}{K}}).$$
The multigraded components of the group $`H^{m_2+m_3+n_2+n_31}(𝒵_{\stackrel{~}{K}})`$ different from that determined by the full subcomplex $`\stackrel{~}{K}_{V_2V_3}`$ do not affect the nontriviality of the Massey product, while the multigraded component corresponding to $`\stackrel{~}{K}_{V_2V_3}`$ is zero since this subcomplex is contractible. The group $`H^{m_1+m_2+n_1+n_21}(𝒵_{\stackrel{~}{K}})`$ is treated similarly. It follows that the Massey product contains a unique nonzero element in its multigraded component and so is nontrivial. ∎
As is well known, the nontriviality of Massey products obstructs formality of manifolds, see e.g. .
###### Corollary 5.6.
For every sphere triangulation $`\stackrel{~}{K}`$ obtained from another triangulation by applying two stellar subdivisions as described above, the 2-connected moment-angle manifold $`𝒵_{\stackrel{~}{K}}`$ is nonformal.
In the proof of Theorem 5.5 the nontriviality of the Massey product is established geometrically. A parallel argument may be carried out algebraically in terms of the algebra $`R^{}(K)`$, as illustrated in the following example.
###### Example 5.7.
Consider the simple polytope $`P^3`$ shown on Figure 4.
This polytope is obtained by cutting two non-adjacent edges off a cube and has 8 facets. The dual triangulation $`K_P`$ is obtained from an octahedron by applying stellar subdivisions at two non-adjacent edges. The face ring is
$$[K_P]=[v_1,\mathrm{},v_6,w_1,w_2]/_{K_P},$$
where $`v_i`$, $`i=1,\mathrm{},6`$, are the generators coming from the facets of the cube and $`w_1,w_2`$ are the generators corresponding to the two new facets, see Figure 4, and
$$_P=(v_1v_2,v_3v_4,v_5v_6,w_1w_2,v_1v_3,v_4v_5,w_1v_3,w_1v_6,w_2v_2,w_2v_4).$$
The corresponding algebra $`R^{}(K_P)`$ has additional generators $`u_1,\mathrm{},u_6,t_1,t_2`$ of total degree 1 satisfying $`du_i=v_i`$ and $`dt_i=w_i`$. Consider the cocycles
$$a=v_1u_2,b=v_3u_4,c=v_5u_6$$
and the corresponding cohomology classes $`\alpha ,\beta ,\gamma H^{1,4}[R^{}(K)]`$. The equations
$$ab=de,bc=df$$
have a solution $`e=0`$, $`f=v_5u_3u_4u_6`$, so the triple Massey product $`\alpha ,\beta ,\gamma H^{4,12}[R^{}(K)]`$ is defined. This Massey product is nontrivial by Theorem 5.5. The cocycle
$$af+ec=v_1v_5u_2u_3u_4u_6$$
represents a nontrivial cohomology class $`[v_1v_5u_2u_3u_4u_6]\alpha ,\beta ,\gamma `$ and so the algebra $`R^{}(K_P)`$ and the manifold $`𝒵_{K_P}`$ are not formal.
In view of Theorem 5.5, the question arises of describing the class of simplicial complexes $`K`$ for which the algebra $`R^{}(K)`$ (equivalently, the Koszul algebra $`[\mathrm{\Lambda }[u_1,\mathrm{},u_m][K],d]`$ or the space $`𝒵_K`$) is formal (in particular, does not contain nontrivial Massey products). For example, a direct calculation shows that this is the case if $`K`$ is the boundary of a polygon.
### 5.4. Toral rank conjecture
Here we relate our cohomological calculations with moment-angle complexes to an interesting conjecture in the theory of transformation groups. This ‘toral rank conjecture’ has strong links with rational homotopy theory, as described in . Therefore this last subsection, although not containing new results, aims at encouraging rational homotopy theorists to turn their attention to combinatorial commutative algebra of simplicial complexes.
A torus action on a space $`X`$ is called *almost free* if all isotropy subgroups are finite. The *toral rank* of $`X`$, denoted $`\mathrm{trk}(X)`$, is the largest $`k`$ for which there exists an almost free $`T^k`$-action on $`X`$.
The *toral rank conjecture* of Halperin suggests that
$$dimH^{}(X;)2^{\mathrm{trk}(X)}$$
for any finite dimensional space $`X`$. Equality is achieved, for example, if $`X=T^k`$.
Moment angle complexes provide a wide class of almost free torus actions:
###### Theorem 5.8 (Davis–Januszkiewicz \[11, 7.1\]).
Let $`K`$ be an $`(n1)`$-dimensional simplicial complex with $`m`$ vertices. Then $`\mathrm{trk}𝒵_Kmn`$.
###### Proof.
Choose an lsop in $`t_1,\mathrm{},t_n`$ in $`[K]`$ according to Lemma 2.3 and write
$$t_i=\lambda _{i1}v_1+\mathrm{}+\lambda _{im}v_m,i=1,\mathrm{},n.$$
Then the matrix $`\mathrm{\Lambda }=(\lambda _{ij})`$ defines a linear map $`\lambda :^m^n`$. Changing $`\lambda `$ to $`k\lambda `$ for a sufficiently large $`k`$ if necessary, we may assume that $`\lambda `$ is induced by a map $`^m^n`$, which we also denote by $`\lambda `$. It follows from Lemma 3.11 that for every simplex $`\sigma K`$ the restriction $`\lambda |_^\sigma :^\sigma ^n`$ of the map $`\lambda `$ to the coordinate subspace $`^\sigma ^m`$ is injective.
Denote by $`T_\mathrm{\Lambda }`$ the subgroup in $`T^m`$ corresponding to the kernel of the map $`\lambda :^m^n`$. Then $`T_\mathrm{\Lambda }`$ is a product of an $`(mn)`$-dimensional torus $`N`$ and a finite group. The intersection of the torus $`N`$ with the coordinate subgroup $`T^\sigma T^m`$ is a finite subgroup. Since the isotropy subgroups of the $`T^m`$-action on $`𝒵_K`$ are of the form $`T^\sigma `$ (see the proof of Theorem 3.12), the torus $`N`$ acts on $`𝒵_K`$ almost freely. ∎
Note that by construction the space $`𝒵_K`$ is 2-connected.
In view of Theorem 5.1, we get the following reformulation of the toral rank conjecture for $`𝒵_K`$:
$$dim\underset{\omega [m]}{}\stackrel{~}{H}^{}(K_\omega ;)2^{mn}$$
for any simplicial complex $`K^{n1}`$ on $`m`$ vertices.
###### Example 5.9.
Let $`K`$ the boundary of an $`m`$-gon. Then the calculation of \[8, Exam. 7.22\] shows that
$$dimH^{}(𝒵_K)=(m4)2^{m2}+42^{m2}.$$
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# The Dark Halo of NGC 5963 as a Constraint on Dark Matter Self-Interaction at the Low Velocity Regime
## 1 Introduction
Numerical studies of structure formation with the collisionless cold dark matter (CDM) scenario predict dark halos with steep central cusps (e.g., Navarro, Frenk & White 1996, 1997, hereafter NFW), whereas most of the rotation curves of dwarf galaxies and low surface brightness (LSB) galaxies suggest that their halos have constant density cores (e.g., Marchesini et al. 2002; de Blok & Bosma 2002, and references therein). Self-interacting dark matter (SIDM) with cross sections per unit of mass in the range $`0.5`$$`6`$ cm<sup>2</sup> g<sup>-1</sup> was proposed as a possible route to reduce the density cusp, as thermal conduction replaces the central cusp by a soft core (Spergel & Steinhardt 2000). A tight constraint on the cross section just under $`0.1`$ cm<sup>2</sup> g<sup>-1</sup> has been derived by considering the formation of giant cluster arcs (Meneghetti et al. 2001) and the size of the cores of galaxy clusters (e.g., Arabadjis, Bautz & Garmire 2002; Lewis, Buote & Stocke 2003; Arabadjis & Bautz 2004). Since this cross section is too small to produce galactic cores, it has been pointed out that a velocity dependent cross section might reduce the effects of self-interaction on cluster scales (e.g., Firmani et al. 2000). For a cross section per unit of mass, which varies as some power of the relative velocity between colliding particles $`\sigma _{\mathrm{dm}}=\sigma _{}(v_0/v_{\mathrm{rel}})^a`$, where $`v_0=(100\mathrm{km}\mathrm{s}^1)`$, there exists a certain range of values of $`\sigma _{}`$ and $`a`$ for which the halos of dwarf galaxies should present a core, while halos of galaxy clusters would not change their core sizes, ellipticities and arcs significantly.
Constraints on the parameters $`\sigma _{}`$ and $`a`$ were obtained by requiring that dwarf galaxies observed today have yet to undergo core collapse and that dark halos must survive the heating from hot cluster halos (Gnedin & Ostriker 2001; Hennawi & Ostriker 2002). All these requirements are satisfied for a very narrow range of parameters $`\sigma _{}=0.5`$$`1`$ cm<sup>2</sup> g<sup>-1</sup> and $`a=0.5`$$`1`$. For cross sections within this suitable range, flat cores with densities of $`0.02`$ M pc<sup>-3</sup> are formed in the central regions of galactic halos as it was confirmed numerically by Colín et al. (2002) in cosmological simulations of SIDM.
Low surface brightness galaxies are ideal to put upper limits on the strength of dark matter (DM) self-interaction at the low-velocity regime (i.e. relative velocities of $`150`$ km s<sup>-1</sup>). In contrast to relatively recently formed objects, like clusters of galaxies, that present unrelaxed mass distributions, LSB galaxies with high central densities may have the highest formation redshifts and, therefore, they have had almost a Hubble time to soften their cusps. In the present study we concentrate on the implications of the halo of the LSB galaxy NGC 5963 for dark matter self-interaction. The high central density and the small core radius of the halo of NGC 5963 suggest that either the DM cross section, at the low velocity regime, is rather small $`0.1`$ cm<sup>2</sup> g<sup>-1</sup> or the dark halo is undergoing an undesirable, dramatic core collapse. Therefore, it seems very unlikely that collisional scattering between DM particles is the main agent for the formation of cores.
We will start in §2 with a description of the basic physics behind the scenario of SIDM and some predictions, which will be used in the remainder of the paper. In §3 we briefly describe the properties of NGC 5963 and present mass models over a range of mass-to-light ratios of the stellar disk. We then compare with predictions of cosmological SIDM simulations to constrain the self-interaction cross section of DM particles (§4). Some implications of the results are discussed in §5.
## 2 Self-Interacting Dark Matter and the Size of the Cores
The mass-density profiles of numerically simulated collisionless CDM halos are commonly parameterized with the analytical Navarro, Frenk & White profile:
$$\rho (r)=\frac{\rho _i}{\left(r/r_s\right)\left(1+r/r_s\right)^2},$$
(1)
where $`r_s`$ is the characteristic radius of the halo and $`\rho _i`$ is related to the density of the Universe at the time of collapse. Within the core region, the dark matter density increases as a power law $`\rho r^1`$ and the velocity dispersion decreases toward the center at the radius $`r_s`$. The parameters $`(\rho _i,r_s)`$ are usually expressed in a slightly different form by the concentration parameter $`c=r_{200}/r_s`$ and the rotation velocity $`V_{200}`$ at radius $`r_{200}`$. The latter is the radius inside which the average overdensity is $`200`$ times the critical density of the Universe.
If dark matter particles experience self-interaction, the scattering thermalizes the inner regions of dark halos, producing a constant-density core (e.g., Burkert 2000). It is expected that the core radius will be comparable to the inversion radius, i.e. the radius at which the velocity dispersion peaks. After reaching the maximum core size, the core can begin to recollapse due to the gravothermal catastrophe. In order to have a core, the relaxation time should be less than $`1/10`$ of Hubble time, so that it should have sufficient time to form a core by thermalization, while the lifetime of its core should exceed $`10`$ Hubble times, so the core should be far from collapse.
In the monolithic scenario in which the halos evolve at isolation, the core radius increases in time as the result of self-interaction between halo particles. Starting with a profile with inversion radius $`r_i`$, the core radius of the pseudoisothermal profile, $`r_c`$, achieves a radius $`0.4r_i`$ in one core-radius relaxation timescale $`t_{\mathrm{rc}}=(3\rho \sigma _{\mathrm{dm}}\stackrel{~}{v})^1`$, with $`\rho `$ the density and $`\stackrel{~}{v}`$ the velocity dispersion, both evaluated at $`0.6r_i`$ (e.g., Kochanek & White 2000)<sup>1</sup><sup>1</sup>1Kochanek & White (2000) estimated the core radius by the point at which the density drops to $`1/4`$ the central density. To estimate $`r_c`$, we have used that $`r_c`$ is related to $`r_{1/4}`$ by the relation $`r_c=r_{1/4}/1.7`$.. The maximum radius that the core reaches before recollapse may be even larger. In fact, for the particular case of a Hernquist (1990) model with break radius $`r_H`$, Kochanek & White (2000) found that $`r_c`$ reaches a maximum value of $`0.23r_H`$. Since the velocity dispersion for a Hernquist profile peaks at $`0.33r_H`$, the core radius $`r_c`$ at the maximum expansion is approximately $`0.7r_i`$.
The assumptions behind the monolithic scenario that the initial configuration follows a NFW profile and that halos evolve at isolation appear to be inconsistent with the SIDM hypothesis in a cosmological context. Cosmological simulations of the formation of SIDM halos, either with a constant cross section or dependent on velocity, have been carried out by Yoshida et al. (2000), Davé et al. (2001), Colín et al. (2002) and D’Onghia, Firmani & Chincarini (2003). The inclusion of dynamically hot material that continuously is accreted onto the halo (in a cosmological context) can prevent efficient heat transfer from the core to the halo. Moreover, it turned out that since collisions already modify halo profiles in the dynamical process of halo growth, there is a smooth trend of increasing core radius with the effective value of $`\sigma _{\mathrm{dm}}`$. A core with a radius $`0.4r_s`$ is formed if the mean number of collisions per particle and per Hubble time in the center halo is about $`4`$$`6`$. We will use this result in §4.
The existence of galaxies with large DM concentrations and central DM densities implies upper limits on the number of collisions per particle and hence on $`\sigma _{}`$. In the following we consider the implications of the concentrated halo in the galaxy NGC 5963.
## 3 The Galaxy NGC 5963 and its Dark Halo
### 3.1 NGC 5963 and its Rotation Curve
NGC 5963 is a relatively isolated LSB galaxy associated with the NGC 5866 group, with a suitable inclination angle of $`45\pm 4\mathrm{deg}`$. Its adopted distance is $`13`$ Mpc, based on the recession velocity derived by Bosma, van der Hulst & Athanassoula (1988) and a Hubble constant of $`70`$ km s<sup>-1</sup> Mpc<sup>-1</sup>. The uncertainty in the distance is probably of order $`3`$ Mpc. The distributions of light and mass were first reported by Romanishin et al. (1982), who found a luminosity profile that exhibits a lens component, i.e. a component with a plateau and a steep outer edge, and a small bulge (see also Simon et al. 2004). The extent of the lens is about $`1.4`$ kpc. The H i rotation curve was constructed by Bosma et al. (1988). These authors emphasized that the dark halo of this galaxy is more concentrated than in normal Sc’s with similar rotation curves. In this sense, NGC 5963 obeys the rule pointed out by Sancisi (2003) that an excess of light corresponds an excess of rotation and viceversa. Beyond a radius of $`7`$ kpc the rotation curve becomes fairly flat with a circular velocity $`v_c`$ $`131`$ km s<sup>-1</sup> (see Fig. 1). If the dark halo can be approximated by the pseudoisothermal sphere, the one-dimensional velocity dispersion for halo particles for NGC 5963 would be of $`\stackrel{~}{v}v_\mathrm{c}/\sqrt{2}=90`$ km s<sup>-1</sup> and, therefore, the halo particles undergo collisions in the low-velocity regime.
In order to study the cuspiness of the mass density profile in the galaxy centers it is required a rotation curve with a high spatial and velocity resolution and a large extension. This can be achieved by combining CO, H$`\alpha `$ data, and radio H i data. NGC 5963 is one of the targets of the sample of galaxies with high-resolution two-dimensional H$`\alpha `$ and CO velocity fields, which are free of beam-smearing, in Simon et al. (2004). Two-dimensional velocity fields are useful in order to avoid systematic uncertainties and to determine the existence of radial motions. Multiple wavelengths studies of the velocity fields may also help to elucidate the origin of systematic errors that can make density profiles appear artificially shallow. The observations in CO and H$`\alpha `$ provide the rotation curve within a galactocentric radius of $`4`$ kpc, whereas the H i rotation curve allows us to trace the potential out to $`10`$ kpc. This galaxy lacks of measurable non-circular motions, i.e. radial velocities $`<5`$ km s<sup>-1</sup>, except perhaps in the rings between $`12^{\prime \prime }`$ to $`30^{\prime \prime }`$, those being not large enough to affect the derivation of the DM density profile. Simon et al. (2004) compared the H$`\alpha `$ and CO velocities at every point and found an excellent agreement between them, with a mean offset of less than $`1`$ km s<sup>-1</sup>, and a scatter of $`7.8`$ km s<sup>-1</sup>. Though the uncertainties in the rotation curve can be reduced by combining both velocity data, we prefer to use the CO rotation curve and its associated errorbars rather than the formal errorbars of the H$`\alpha `$ rotation curves, which simply indicate that the Gaussian fit to the profile was well-determined. In particular, the CO rotation curve is less sensitive to asymmetric drift corrections and other magnetic terms, due to the smaller small-scale motions of the molecular clouds compared to the ionized gas (Sánchez-Salcedo & Reyes-Ruiz 2004). In addition, other uncertainties in mass models, as the mass-to-light ratio of the different components, impose a minimum error in any fit. Figure 1 shows the CO rotation curve in the inner galaxy and the H i rotation curve in the outer parts, with error bars indicating the uncertainties in the observations.
The present galaxy is an ideal laboratory with which to test the effects of collisions between DM halo particles in the low velocity regime, not only for the properties of the galaxy but also for the quality of the data.
### 3.2 Mass Models and Adiabatic Contraction
We have constructed different mass models for this galaxy. First we calculate the rotation curve due to the neutral (H i$`+`$He) gas mass from the radial distribution of the H i surface density from Bosma et al. (1988). The surface brightness profiles in the optical and near-IR bands given by Simon et al. (2004) were used to derive the disk mass distribution by adopting a constant mass-to-light ratio for the disk. The resulting mass distribution was converted into a disk rotation curve. The difference between the sum of the squares of the observed velocity and the stellar plus gaseous rotation velocities was converted into a halo mass distribution after assuming a certain form for the halo. Along this paper, all the fits to the CO rotation curve were performed by a $`\chi ^2`$-minimization. To make the best use of both the CO and H i data, we have used a hybrid rotation curve. It consists of the CO data over the range of radii where available and $`21`$ cm data to define the outermost points. This was accomplished by minimizing $`\chi ^2`$ between $`0`$ and $`4`$ kpc but requiring that between $`5`$ and $`10`$ kpc, the goodness of the fit must be better than $`1\sigma `$ level, while rejected otherwise. In this way, the fits are statistically consistent with the H i rotation curve at large radii.
Different mass decompositions were fitted for the stellar mass-to-light ratio in the R-band varying from $`0`$ (minimum disk) up to $`1.2`$ M/L$`{}_{}{}^{R}{}_{}{}^{}`$ (maximum disk). We will denote by $`\mathrm{{\rm Y}}_{}`$ the stellar mass-to-light ratio in the R-band in solar units. It is noteworthy that $`\mathrm{{\rm Y}}_{}=1.2`$ is the maximum value allowed by the smooth rotation curve but avoiding a hollow halo in the inner two points. However, it is submaximal in the sense that the maximum contribution of the stellar disk to the rotation curve is less than $`0.8`$ times the observed maximum rotation. In practice, the stellar curve is scaled until the inner points match those of the smooth curve. Depending on how many points are used, the maximum disk mass-to-light ratios may vary as much as $`15\%`$. In the case of NGC 5963 and using the inner five points, Simon et al. (2004) obtained a dynamical value for the maximum mass-to-light ratio larger than the value quoted above, but it appeared unrealistic and inconsistent with the predictions from the galaxy colours by the Bell et al. (2003) population synthesis models. The dynamical maximum disk $`\mathrm{{\rm Y}}_{}`$ inferred by using two points is fully consistent with these predictions.
Fig. 2 shows the best-fitting mass models for different values of $`\mathrm{{\rm Y}}_{}`$ with both the NFW halo and the cored pseudoisothermal halo:
$$\rho ^{}(r)=\frac{\rho _0^{}r_\mathrm{c}^2}{r^2+r_\mathrm{c}^2}.$$
(2)
Primed quantities are used where appropriate to denote the (observed) final states after adiabatic compression of the dark halo by the infall of the baryons as they cool and settle into a disk. The parameters of each model, the reduced $`\chi ^2`$ of the fit and the probability $`p`$ that the data and the model could result from the same parent distribution are given in Table 1. It is important to remark that the observation that pseudoisothermal fits produce low reduced $`\chi ^2`$ values does not demonstrate the presence of a constant-density core. Simon et al. (2004) have shown that a power-law density profile $`\rho r^{1.2}`$ fits also very well the rotation curve of NGC 5963.
In order to compare with simulations which do not include the baryonic physics, we need an estimate of the DM profile before the adiabatic compression by the baryon condensation, i.e. the unprimed parameters of the halo. A cored profile is expected after adiabatic decompression. Since the core baryons cause the dark profile to steepen, we have fitted the observed rotation curve, for a given $`\mathrm{{\rm Y}}_{}0.2`$, with a cuspy profile (for convenience we have chosen the NFW profile) and recovered the halo parameters $`(\rho _0,r_c)`$ of the approximately pseudoisothermal sphere before compression, following the standard model (e.g., Blumenthal et al. 1986; Flores et al. 1993). For $`\mathrm{{\rm Y}}_{}<0.2`$, the effect of compression by the baryons is so small that a cored profile is assumed before and after contraction. The recipe of Blumenthal et al. (1986) assumes that the matter distribution is spherically symmetric and that particles move on circular orbits. Jesseit et al. (2002) verified that if the disk formation is smooth, the contracted halos have density profiles that are in excellent agreement with the standard predictions. If the disk formation involves clumpy, cold streams or if the bulk of the central stars formed when the halo mass was still being assembled, the level of halo contraction may be smaller. Gnedin et al. (2004) suggested an improvement of the standard model of the response of the halo to condensation of baryons by a simple modification of the assumed invariant (see Gnedin et al. 2004 for details). They show that the standard model systematically overpredicts the contraction. Moreover, there exist other empirical evidence which suggests that adiabatic contraction is avoided (e.g., Loeb & Peebles 2003; Dutton et al. 2005). Therefore, our procedure provides an upper limit on the effect of adiabatic contraction.
For illustration, Fig. 3 shows the contribution to the rotation curve of the dark halo after correcting by adiabatic contraction, in the most extreme case of maximum disk $`\mathrm{{\rm Y}}_{}=1.2`$. The fit to it adopting the pseudoisothermal profile is also shown. Note that if we demand that the density of the halo must decrease monotonically with increasing radius, values $`\mathrm{{\rm Y}}_{}>1.2`$ would not be permitted in the standard adiabatic decompression approach.
In Table 1, best-fitting halo parameters are reported as a function of $`\mathrm{{\rm Y}}_{}`$. Even adopting the maximum disk solution, the halo of NGC 5963, at present time, is extraordinary concentrated, with a central dark matter density greater than $`0.35`$ M pc<sup>-3</sup> and a core radius $`<1`$ kpc. A thorough discussion aimed to identify what about NGC 5963 makes it unique was given by Simon et al. (2004) and need not be repeated here. In the next sections, we will discuss the implications of the halo of NGC 5963 for the self-interaction of DM particles.
## 4 Constraints on the Cross-Section of DM Particles at the Low Velocity Regime
In principle we do not really know whether the core of the galaxy under consideration is still in the expansion phase or undergoing gravothermal core collapse aided by the gravitational contraction due to the baryons. In fact, the core contraction of the halo by the collapsing baryons may trigger a fast core collapse (e.g., Kochanek & White 2000, and §4.3). This effect may be important for $`\mathrm{{\rm Y}}_{}`$ ratios close to the maximum disk value. Using cosmological simulations of SIDM as a calibrator, we first put limits on the cross section of DM particles assuming that the core of NGC 5963 is in the expanding phase due to heat transfers inwards (§4.1). For large values of $`\mathrm{{\rm Y}}_{}`$, however, core collapse triggered by baryon condensation may have dramatic consequences for the longevity of the core. We have to demand that the halos are in no danger of collapsing. This requirement implies another bound on the cross section, which is discussed in §4.3.
The galaxy halo profiles and the tight cross section limits coming from clusters of galaxies can be reconciled only if the cross section were inversely proportional to the halo velocity dispersion. Nevertheless, except for subhalos undergoing the action of the hot halo environment, most of the relevant results found in simulations with a cross section independent of halo velocity dispersion are valid if the value of the cross section is interpreted as that appropriate to the velocity dispersion of the halo being considered. Therefore, since we are primarily interested in a single object, we can gain additional physical insight by making also use of simulations with a velocity-independent cross section.
### 4.1 Halos in Midsized Galaxies under SIDM: Calibrating Collisional Effects
In cosmological simulations of structure formation in SIDM scenario with effective cross sections $`0.5`$$`1`$ cm<sup>2</sup> g<sup>-1</sup>, midsized halos (masses $`10^{1011}`$ M) present cores with radii in the range $`r_c2.5`$$`5`$ kpc and central densities $`\rho _00.01`$$`0.06`$ M pc<sup>-3</sup> (Yoshida et al. 2000; Davé et al. 2001; Colín et al. 2002; D’Onghia et al. 2003). For instance, D’Onghia et al. (2003) derive $`r_c=5`$$`6`$ kpc and $`\rho _00.02`$ M pc<sup>-3</sup> for a halo with maximum circular velocity of $`120`$ km s<sup>-1</sup>. Therefore, effective cross sections in the range $`0.5`$$`1`$ cm<sup>2</sup> g<sup>-1</sup> are suitable to reproduce soft cores in late-type galaxies<sup>2</sup><sup>2</sup>2In order to compare simulations, we notice that if the cross section is assumed to be inversely proportional to the collision velocity, $`\sigma _{\mathrm{dm}}=\sigma _{}(v_0/v_{\mathrm{rel}})`$, the effective cross section for a halo with (one-dimensional) velocity dispersion $`\stackrel{~}{v}`$ is $`\sigma _{}/(\sqrt{\pi }\stackrel{~}{v}_{100})`$, where $`\stackrel{~}{v}_{100}=(\stackrel{~}{v}/100\mathrm{km}\mathrm{s}^1)`$. To show this relation one has to remind that, for a Maxwellian velocity distribution, $`1/v_{\mathrm{rel}}=1/(\sqrt{\pi }\stackrel{~}{v})`$..
All the simulations with $`\sigma _{\mathrm{dm}}=\sigma _{}(v_0/\stackrel{~}{v})`$ also shown that for $`\sigma _{}=0.5`$$`1`$ cm<sup>2</sup> g<sup>-1</sup>, the number of collisions per particle per Hubble time at the center of the halos, was between $`4`$$`6`$, roughly independent of the halo mass. In fact, while a few collisions (between $`2`$ and $`3`$) at the halo center are enough to produce a constant density core (Yoshida et al. 2000), $`4`$$`6`$ collisions are sufficient to produce a core with $`r_c0.4r_{s0}`$, with $`r_{s0}`$ the scale radius of the NFW halo when it is resimulated in the standard, collisionless case ($`\sigma _{\mathrm{dm}}=0`$). The reason is that each scattering produces a change $`\mathrm{\Delta }vv`$, and particles escape the core in a single scattering.
Rescaling the number of collisions per particle at the center of the cluster labelled S1Wa in Yoshida et al. (2000), we find
$$N_{\mathrm{col}}3\left(\frac{\rho _0}{0.02\mathrm{M}_{}\mathrm{pc}^3}\right)\left(\frac{\stackrel{~}{v}}{100\mathrm{km}\mathrm{s}^1}\right)\left(\frac{\sigma _{\mathrm{dm}}}{1\mathrm{cm}^2\mathrm{g}^1}\right).$$
(3)
We stress here that $`\rho _0`$ is the final unprimed central density of the core, after thermalization. Let us estimate $`N_{\mathrm{col}}`$ in the halo center of NGC 5963 for mass models with $`\mathrm{{\rm Y}}_{}<0.7`$. For a DM cross section $`\sigma _{\mathrm{dm}}\stackrel{~}{v}_{100}`$ between $`0.5`$ and $`1.0`$ cm<sup>2</sup> g<sup>-1</sup>, where $`\stackrel{~}{v}_{100}=(\stackrel{~}{v}/100\mathrm{km}\mathrm{s}^1)`$, and a central density $`\rho _00.4`$ M pc<sup>-3</sup>, the mean number of collisions is $`N_{\mathrm{col}}20`$$`45`$. Hence, we expect NGC 5963 having a core radius greater than $`0.4r_{s0}`$. In fact, from the scaling laws of Yoshida et al. (2000), we see that the core radius reaches a size $`1.5r_{s0}`$ when the collision rate per particle at halo center is $`30`$ per Hubble time. In order to quantify the expected size of the core of this galaxy under SIDM, we need its characteristic radius $`r_{s0}`$.
An estimate of $`r_{s0}`$ can be inferred for NGC 5963 by fitting the rotation curve beyond $`5`$ kpc with a NFW profile. The observed inner rotation curve cannot be included as it is altered as a consequence of DM self-interactions and the adiabatic contraction that causes the DM profile to steepen. Since there are many combinations of the parameters $`(c,V_{200})`$ that can fit the rotation curve, we need to choose either $`c`$ or $`V_{200}`$. Expected values of the concentration in $`\mathrm{\Lambda }`$CDM cosmology have a $`2\sigma `$ range from $`5`$ to $`22`$ (Eke et al. 2001; Jing & Suto 2002). We will take a value $`c=20`$, which lies among the largest values in the $`2\sigma `$ uncertainty. The exact value of $`r_{s0}`$ also depends slightly on $`\mathrm{{\rm Y}}_{}`$. For the intermediate case $`\mathrm{{\rm Y}}_{}=0.6`$ (and $`c=20`$), we find that the characteristic radius is $`r_{s0}6`$ kpc. For the average value predicted by simulations, $`c=13`$, we found $`r_{s0}=11.7`$ kpc. Hence, as a conservative calibration, we will assume that for $`N_{\mathrm{col}}6`$, a core radius $`r_c0.4r_{s0}2.5`$ kpc should have been formed.
In Table 1 we see that for models with $`\mathrm{{\rm Y}}_{}<0.7`$, the inferred core radius before contraction in NGC 5963 is smaller than $`1.2`$ kpc; the question that arises is how many collisions suffice to develop a core with $`r_c1`$ kpc. To answer this question we have examined the SIDM simulations of the evolution of a Hernquist halo with total mass $`M_T`$ and break radius $`r_H`$ in Kochanek & White (2000). These authors carried out simulations surveying the dimensionless cross section $`\widehat{\sigma }_{\mathrm{dm}}=M_T\sigma _{\mathrm{dm}}/r_H^2`$. According to their Fig. 2c, the core density in the phase of core expansion, and for halos with different $`\widehat{\sigma }_{\mathrm{dm}}`$, goes as $`\widehat{\sigma }_{\mathrm{dm}}^\eta `$, with $`\eta <1/2`$. Moreover, at the time that the core of the simulation with dimensionless cross section $`\widehat{\sigma }_{\mathrm{dm}}=3`$ reaches its maximum radius, the run $`\widehat{\sigma }_{\mathrm{dm}}=0.3`$ has already developed a core of half this radius, with a mean collision count at halo center $`5`$ times smaller. Applying this scaling to NGC 5963 it holds that a mean of $`6/5`$ collisions per particle at the center, will suffice to produce a core of $`2.5/2`$ kpc. Although these estimates are based on simulations of the relaxation of an isolated galaxy having initially a cuspy Hernquist profile, it was also found in cosmological simulations that a few collisions (exceeding $`2`$) suffice to develop a kpc-size core (e.g., Yoshida et al. 2000). To keeps matter simple we will adopt the generous condition that only if $`N_{\mathrm{col}}2`$ then $`r_c1`$ kpc.
Putting together, we take the following, rather conservative, relationships as a calibrator of the collisional effects:
$$\mathrm{if}N_{\mathrm{col}}=6r_c=2.5\mathrm{kpc},$$
(4)
$$\mathrm{if}N_{\mathrm{col}}=2r_c1.0\mathrm{kpc}.$$
(5)
Suppose that $`\sigma _{\mathrm{dm}}\stackrel{~}{v}^1`$. The generalization for a power-law $`\sigma _{\mathrm{dm}}\stackrel{~}{v}^a`$ is obvious and it will be ignored in the interest of simplicity. To place an upper limit on $`\sigma _{\mathrm{dm}}\stackrel{~}{v}_{100}`$, we proceed as follows. We first derive the best-value fitting parameters ($`r_c,\rho _0)`$ for a given $`\mathrm{{\rm Y}}_{}`$. If $`1`$ kpc $`<r_c<3`$ kpc, we estimate $`N_{\mathrm{col}}`$ by interpolating Eqs. (4)-(5) linearly. Once $`N_{\mathrm{col}}`$ is known, Equation (3) immediately provides us with an upper limit on $`\sigma _0\sigma _{\mathrm{dm}}\stackrel{~}{v}_{100}`$. If $`r_c`$ is smaller than $`1`$ kpc, we will assume that $`N_{\mathrm{col}}2`$, which gives:
$$\sigma _00.7\mathrm{cm}^2\mathrm{g}^1\left(\frac{\rho _0}{0.02\mathrm{M}_{}\mathrm{pc}^3}\right)^1.$$
(6)
For instance, if we have a situation in which $`r_c<1`$ kpc and $`\rho _0=0.4`$ M pc<sup>-3</sup>, we would obtain $`\sigma _00.035`$ cm<sup>2</sup> g<sup>-1</sup>.
### 4.2 Results
Figure 4 shows the resulting upper limits on $`\sigma _0`$ versus $`\mathrm{{\rm Y}}_{}`$, following the procedure described in the previous subsection. Confidence levels were calculated as a measure of the sensitivity of the results to the goodness of the fit to the rotation curve. Other sources of uncertainties in the parameters of the model, as the distance to the galaxy, its inclination, or the adopted characteristic radius $`r_{s0}`$, were not taken into account in these contours. We must warn that these error bands cannot be considered as real probability indicators because the velocities and their errors are not free of systematic effects.
Our derivation overestimates the upper limit on $`\sigma _0`$ because we have assumed in the derivation of the halo parameters that baryons only act to hasten core contraction. Nevertheless, other phenomena as galactic bars, outflows, massive black holes in the halos, and other processes associated with AGN or star formation have been proposed as mechanisms to erase cuspy halos. We are being conservative in ignoring these phenomena because incorporating their effects would only serve to place more stringent constraints on $`\sigma _0`$.
The bound on $`\sigma _0`$ depends strongly on $`\mathrm{{\rm Y}}_{}`$. From Fig. 4 we see that models with $`\mathrm{{\rm Y}}_{}0.6`$ are consistent with $`\sigma _0`$ just under $`0.1`$ cm<sup>2</sup> g<sup>-1</sup>. For $`\mathrm{{\rm Y}}_{}=0.35`$, $`\sigma _0>0.05`$ cm<sup>2</sup> g<sup>-1</sup> can be excluded virtually at $`2\sigma `$ confidence interval. Within the range $`0<\mathrm{{\rm Y}}_{}<0.7`$, our limit excludes the interval proposed to explain the flat mass profiles in galaxies.
For $`\mathrm{{\rm Y}}_{}>1.0`$, the present analysis only disfavours the interval of astrophsical interest, $`\sigma _0=0.3`$$`1`$ cm<sup>2</sup> g<sup>-1</sup>, at less than $`68\%`$ confidence level. We should notice, however, that if the contribution of the baryons to the potential well is important, the core may be in danger of undergoing gravothermal collapse. The requirement that the core is not undergoing a fast collapse will provide a more stringent upper limit for the DM cross section for $`\mathrm{{\rm Y}}_{}>0.8`$.
### 4.3 Core Collapse
In the latter section we found that a cross section $`\sigma _00.5`$ cm<sup>2</sup> g<sup>-1</sup> and $`\mathrm{{\rm Y}}_{}>1.0`$ would be marginally consistent with a picture in which the halo of NGC 5963 might have formed a core with a radius $`2`$$`2.5`$ kpc, although significantly reduced because of the adiabatic contraction by baryons. Nevertheless, we must also require that this galaxy has yet to undergo core collapse. For these high $`\mathrm{{\rm Y}}_{}`$, the cooling baryons will compress the dark matter and accelerate the core collapse in different ways (e.g., Kochanek & White 2000). By raising the density due to the contraction, the relaxation time drops. In addition, the adiabatic compression produces a steeper central density cusp which further speeds up the ultimate evolution of the system. Moreover, the initial inversion of the temperature (velocity dispersion of DM particles) profile responsible for the core expansion may be erased by the adiabatic heating of DM particles, leading directly to a core collapse with no room for the expansion phase. We need to estimate the characteristic time for core collapse in the case $`\mathrm{{\rm Y}}_{}0.8`$, i.e. when the effects of contraction are important. In their numerical simulations, Kochanek & White (2000) assumed as initial conditions a Hernquist profile with a cuspy core. In our case, i.e. for values $`\mathrm{{\rm Y}}_{}0.8`$, it is realistic to start with such an initial configuration because a cusp in the density profile is a natural consequence of the adiabatic contraction caused by the baryons. In this regard, we can make use of the quantitative study of Kochanek & White (2000).
The evolution of the collapse depends on the ratio between the collision mean free path, $`\lambda `$, to the local gravitational scale height $`H`$ (Balberg et al. 2002). Systems initially having ratios $`(\lambda /H)_{\mathrm{edge}}10`$ at the outer edge of the core, have a lifetime of only a few relaxation times<sup>3</sup><sup>3</sup>3It is simple to show that $`\widehat{\sigma }_{\mathrm{dm}}`$ and $`(\lambda /H)_{\mathrm{edge}}`$ are related by the relation $`(\lambda /H)_{\mathrm{edge}}2.2/\widehat{\sigma }_{\mathrm{dm}}`$. (Quinlan 1996; Kochanek & White 2000; Balberg et al. 2002). For the range of $`\widehat{\sigma }_{\mathrm{dm}}`$ explored in Kochanek & White (2000), the central density increases by an order of magnitude in a timescale $`t_{c,10}4\widehat{\sigma }_{\mathrm{dm}}^{1/2}t_{\mathrm{rc}}`$, being $`\widehat{\sigma }_{\mathrm{dm}}=2\pi \rho _ir_s\sigma _{\mathrm{dm}}`$ for the NFW profile. In particular, for $`\widehat{\sigma }_{\mathrm{dm}}=0.3`$, this timescale is only $`2.2`$ times the core-radius relaxation time. This result is troublesome for the survival of the core because it recollapses quickly after its formation and could not persist until today (see §2).
Let us calculate the core timescale $`t_{c,10}`$ for the halo of NGC 5963. As said before, because of the adiabatic compression by the baryons, a cuspy halo, such as the NFW profile, is expected even in the presence of DM collisions. The best NFW model for $`\mathrm{{\rm Y}}_{}=1.2`$ corresponds to $`c=15.6`$ and $`V_{200}=109`$ km s<sup>-1</sup>, implying a core-radius timescale:
$$t_{\mathrm{rc}}=\frac{1}{3\rho _i\sigma _{\mathrm{dm}}\stackrel{~}{v}}=0.8\mathrm{Gyr}\left(\frac{\sigma _0}{1\mathrm{cm}^2\mathrm{g}^1}\right)^1.$$
(7)
For the above halo parameters, $`\widehat{\sigma }_{\mathrm{dm}}=2\pi \rho _ir_s\sigma _{\mathrm{dm}}=0.15\stackrel{~}{v}_{100}^1(\sigma _0/1\mathrm{cm}^2\mathrm{g}^1)`$ and hence the central density will be enhanced by a factor $`10`$ in the timescale
$$t_{c,10}4\widehat{\sigma }_{\mathrm{dm}}^{1/2}t_{\mathrm{rc}}=1.25\mathrm{Gyr}\left(\frac{\stackrel{~}{v}}{100\mathrm{km}\mathrm{s}^1}\right)^{1/2}\left(\frac{\sigma _0}{1\mathrm{cm}^2\mathrm{g}^1}\right)^{1/2}.$$
(8)
Taken $`\sigma _0=0.5`$ cm<sup>2</sup> g<sup>-1</sup>, we get $`t_{c,10}2`$ Gyr. This estimate has uncertainties of as much as a factor $`2`$. This short timescale indicates that the dark halo may have sufficient time to increase its central density by a factor $`100`$-$`1000`$. In order to prevent the dark halo from such dramatic evolution, we must demand that $`t_{\mathrm{rc}}`$ one Hubble time, implying $`\sigma _00.08`$ cm<sup>2</sup> g<sup>-1</sup>. Proceeding in the same way we obtain $`\sigma _00.05`$ cm<sup>2</sup> g<sup>-1</sup> but now for the mass model with $`\mathrm{{\rm Y}}_{}=1.0`$. The line connecting these two points in Fig. 4 delimits the region where cross sections are probably small enough to avoid catastrophic core collapse in the halo of NGC 5963. The corresponding line at the $`1\sigma `$ confidence is also plotted. We see that the new constraint, which accounts for the role of baryons when $`\mathrm{{\rm Y}}_{}>0.85`$, is tighter than the one derived in the previous subsections. This all but removes the permitted window for cross sections $`\sigma _0=0.3`$$`1`$ cm<sup>2</sup> g<sup>-1</sup> in mass models with $`\mathrm{{\rm Y}}_{}>0.85`$, that remained open. Combining both constraints, the maximum of the permitted value $`\sigma _00.2`$ cm<sup>2</sup> s<sup>-1</sup>, at $`2\sigma `$ confidence level, occurs for $`\mathrm{{\rm Y}}_{}0.7`$.
## 5 Discussion and conclusions
Apart from the obvious interest for the still unknown nature of dark matter, the possibility of it having a nonzero self-interaction cross section has other astrophysical implications. SIDM was suggested as a route to form cores in LSB and dwarf galaxies. The fact that the required cross sections are comparable to the cross section for particles interacting with each other via the strong force, has led to speculate that DM particles could interact with both themselves and with baryons through the strong force (Wandelt et al. 2001). The interaction of dark matter with protons might contribute to reheat the intracluster medium in the central regions of clusters of galaxies (Qin & Wu 2001; Chuzhoy & Nusser 2004).
A large range of parameters space is being ruled out by current experimental and astrophysical bounds. Previous studies have gradually whittled down the DM cross section allowed to solve the cuspy problem of halos in $`\mathrm{\Lambda }`$CDM cosmology. In addition to statistical studies, the analysis of individual galaxies can give additional constraints on the strength of DM self-interaction cross section.
The dark matter distribution in NGC 5963 is challenging for any model designed to form central cores in dwarf and LSB galaxies. Apparently, its halo distribution seems at odds with the results of SIDM simulations with cross sections $`\sigma _{\mathrm{dm}}\stackrel{~}{v}_{100}=0.3`$$`1`$ cm<sup>2</sup> g<sup>-1</sup>, which produce central densities $`0.02`$ M pc<sup>-3</sup>, fairly independent of the halo mass, and core radii $`2.5`$$`5`$ kpc. The highly concentrated halo of NGC 5963 implies upper limits for the interaction of dark matter particles, unless we change our assumption of constant $`\mathrm{{\rm Y}}_{}`$. A higher $`\mathrm{{\rm Y}}_{}`$ for the disk than for the lens is required to have a less concentrated halo. However, the observed B-V colours of the disk and the lens do not favour this possibility (see Bosma et al. 1988 for a discussion).
Great efforts have been made to constrain the collisional cross section of DM particles. It was suggested that this cross section should be small enough so the core halo of dwarf and LSB galaxies would not collapse in a Hubble time (e.g., Hennawi & Ostriker 2002). Hence, the halo of NGC 5963 by itself cannot be going through the gravothermal contraction phase. However, core collapse may be induced by the action of the baryons as they deepen the potential well and compress the core in an adiabatic process. Therefore, the core of NGC 5963 may be expanding if the thermalization has not been completed, or shrinking if the core is adiabatically compressed, which may occur only at large values of $`\mathrm{{\rm Y}}_{}`$.
Either the core of NGC 5963 is expanding or contracting in size, a tight constraint $`\sigma _0<0.2`$ cm<sup>2</sup> g<sup>-1</sup> at $`95\%`$ confidence level, which indicates that our results are robust to reasonable DM uncertainties, is derived. This upper limit for $`\sigma _0`$ may be overestimated as much as a factor $`2`$ because of our conservative assumptions. Thus, the original motivation for SIDM of lowering the core densities of galactic halos require collisional cross sections too large to be consistent with the halo of NGC 5963.
One of the largest uncertainties in any mass model is the precise value of $`\mathrm{{\rm Y}}_{}`$ because depends on extinction, star formation history, etc. Based on various considerations, e.g., stellar population synthesis, stellar counts and kinematics in the solar neighbourhood and kinematics of external galaxies, some authors argue that $`\mathrm{{\rm Y}}_{}^I2h`$ in the I-band (e.g., Mo & Mao 2000, and references therein) and $`\mathrm{{\rm Y}}_{}^R1.2`$ in the R-band (Simon et al. 2004). This would mean that the constraint derived from the core collapse discussed in §4.3 is the most relevant, and suggests an upper limit $`\sigma _00.1`$ cm<sup>2</sup> g<sup>-1</sup>. This value coincides with the upper limit inferred in clusters of galaxies (e.g., Meneghetti et al. 2001; Arabadjis & Bautz 2004).
Upper limits in the $`0.02`$$`0.1`$ cm<sup>2</sup> g<sup>-1</sup> range were derived by Hennawi & Ostriker (2002) from the mass of supermassive black holes in the centers of galaxies. For effective cross sections $`\sigma _{\mathrm{dm}}>0.02`$ cm<sup>2</sup> g<sup>-1</sup>, the accretion of SIDM onto seed black holes would produce excessively massive black holes. However, there are serious uncertainties associated with this limit because they make use of the hypothesis that initially the density DM distribution follows the NFW profile, whereas the numerical study of Colín et al. (2002) of SIDM cosmology suggests that the NFW profile at the center of the halos is not achieved at any time. Therefore, the otherwise excessive accretion of matter onto the black hole may not occur (Colín et al. 2002).
In order to reconcile the collisional DM hypothesis as a viable explanation of the formation of the constant density cores with the halo of NGC 5963, we should identify potential ways able to compensate the evacuation of DM in the central parts caused by the scattering of DM or to forestall the core collapse, depending whether the core of NGC 5963 is expanding or collapsing. Tidal redistribution of mass at the halo center cannot be efficient in NGC 5963 because the nearest large galaxy in the group is at a projected distance of $`430`$ kpc. One could relax our simplifying assumption that the halo consists of a well mixed homogeneous DM distribution and to consider a clumpy medium. The final density profile will be the result of two competing effects. On the one hand, the mass infall associated with the spiraling of putative massive clumps of DM ($`10^{67}`$ M) towards the galactic center by dynamical friction, produces a replenishment of material and deepens the central potential well. On the other hand, dynamical friction heating may be effective in flattening the inner DM profile. Unlike clusters of galaxies in which substructure is important in determining the final DM distribution (e.g., Nipoti et al. 2004), halos of LSB galaxies are already largely assembled at $`z3`$ and hence, they have had sufficient time to disrupt the halo substructure, forming a smooth DM distribution within the luminous radius, unless NGC 5963 had a prominent mass aggregation history. However, an anomalous accretion history would be very unsatisfactory because it would have dramatic consequences for the disk itself, producing an excessive dynamical heating or even its destruction. Moreover, Ma & Boylan-Kolchin (2004) argued that energy deposition by merging dark matter substructures likely flatten density profiles. In the lack of any of those potential mechanisms capable to make significant mass redistribution, we think that our analysis is robust, even though it is based on the halo of a single galaxy.
The halo of NGC 5963 is problematic for any model whose mechanism to produce large cores in LSB galaxies depends on collisions between DM particles (e.g., annihilating dark matter), reducing the parameters space and suggesting new directions for DM search. Our analysis also places strong constraints to the non-hadronic exotic Q-balls as a dark matter candidate in galaxies. These “particles” might be arranged to have mutual collisions with a large cross section and, in addition, Q-balls can stick together after collision, reducing the self-interaction as scattering proceeds (Kusenko & Steinhardt 2001). In order to be consistent with the halo of NGC 5963, self-interactions between Q-balls should shut off to a negligible value after $`1`$$`2`$ collisions per particle in order for the initial scatterings to smooth out halo cusps but avoiding gravothermal collapse.
Perhaps the solution of the halo core problem resides in decaying dark matter (Cen 2001a,b; Sánchez-Salcedo 2003). If dark matter particles in galactic halos decay to stable particles with a recoiling velocity of a few tens of kilometers per second, then a fraction of LSB galaxies can still present a substantial concentration. This novel dynamics associated with decaying dark matter was illustrated by Sánchez-Salcedo (2003) for the case of the LSB galaxy NGC 3274. We believe that the halos of NGC 5963 and NGC 3274 are not pathological cases as far as their DM distributions concern. There are additional independent evidence on the existence of galaxies with large central densities. Loewenstein & Mushotzky (2002), in an as-yet-unpublished work, have determined the enclosed mass profile for the elliptical galaxy NGC 4636 using X-ray observations. For this galaxy the central density in models with dark matter cores is higher than expected in the SIDM scenario with $`a1`$ (Loewenstein & Mushotzky 2002).
Future observations of LSB galaxies containing a very low density of luminous material even in the inner parts, so that the observed dynamics should be dominated by the gravitational forces of the dark halo at small radii as well as large radii, will be able to determine whether cores are produced by the gravitational interaction between the luminous and dark matter, and will provide further constraints on the nature of DM in galaxies.
I thank Tony García Barreto and Alberto Bolatto for helpful discussions. I am grateful to the anonymous referee for constructive comments. This work was supported by CONACYT project 2002-C40366.
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# On divergence, relative entropy and the substate property
## 1 Introduction
We consider three measures of distinguishability between quantum states and show various relationships between them. The first measure that we consider is called divergence. It was first considered in \[JRS02\] and is defined as follows:
###### Definition 1 (Divergence)
Let $`\rho ,\sigma `$ be two quantum states. Let $`M`$ be a Positive operator-valued measurement($`\mathrm{𝖯𝖮𝖵𝖬}`$). Then divergence between them denoted $`D(\rho |\sigma )`$, is defined as,
$$D(\rho |\sigma )\underset{𝖬:\mathrm{𝖯𝖮𝖵𝖬}}{\mathrm{max}}\text{Tr}\text{ }M\rho \mathrm{log}\frac{\text{Tr}\text{ }M\rho }{\text{Tr}\text{ }M\sigma }$$
The second is the well known measure called the relative entropy also known as Kullberg-Liebeck divergence (\[NC00\]). It is defined as follows:
###### Definition 2 (Relative entropy)
Let $`\rho ,\sigma `$ be quantum states. Then relative entropy between them, denoted $`S(\rho |\sigma )`$ is defined as,
$$S(\rho |\sigma )\text{Tr}\text{ }(\rho \mathrm{log}\rho \rho \mathrm{log}\sigma )$$
The third measure that we consider we call the substate property. It was also first considered in \[JRS02\]. It is defined as follows:
###### Definition 3 (Substate property)
Two states $`\rho `$ and $`\sigma `$ are said to have the $`k`$-substate property if, $`r>1,\rho _r`$ such that
$$\rho \rho _r_t2/\sqrt{r}\text{ and }\sigma (1\frac{1}{r})\frac{\rho _r}{2^{rk}}0$$
## 2 Results in this article
The following theorem is a compilation of all the results in this article.
###### Theorem 1
Let $`\rho `$ and $`\sigma `$ be two quantum states in $`^n`$. Then,
1. $`D(\rho |\sigma )S(\rho |\sigma )+1`$. This is not new and was shown in \[JRS02\]. We present a proof for completeness.
2. Given classical distributions $`P`$ and $`Q`$ on $`[n]`$, $`S(P|Q)D(P|Q)(n1).`$
3. $`S(\rho |\sigma )D(\rho |\sigma )(n1)+\mathrm{log}n`$.
4. There exists classical distributions $`P`$ and $`Q`$ on $`[n]`$ such that,
$$S(P|Q)>(D(P|Q)/21)(n2)1.$$
5. (Substate theorem) $`\rho `$ and $`\sigma `$ have the $`(8D(\rho |\sigma )+14)`$-substate property. This is not new and was shown in \[JRS02\]. Please refer to \[JRS02\] for a proof.
6. (Converse of substate theorem for classical distributions) If distributions $`P`$ and $`Q`$ have the $`k`$-substate property then $`D(P|Q)2k+2`$.
7. If the following strong $`k`$-substate property holds, i.e. $`\sigma \frac{\rho }{2^k}0`$, then $`S(\rho |\sigma )k`$.
8. There exists a POVM such that, if $`P`$ and $`Q`$ are resulting classical distributions then,
$$S(P|Q)\frac{S(\rho |\sigma )\mathrm{log}n}{n1}1.$$
Proof:
1. Let $`M`$ be the POVM that achieves $`D(\rho |\sigma )`$. Let $`\text{Tr }M\rho \stackrel{\mathrm{\Delta }}{=}p`$ and $`\text{Tr }M\sigma \stackrel{\mathrm{\Delta }}{=}q`$.
$`S(\rho |\sigma )`$ $``$ $`p\mathrm{log}{\displaystyle \frac{p}{q}}+(1p)\mathrm{log}{\displaystyle \frac{(1p)}{(1q)}}`$
$`>`$ $`p\mathrm{log}{\displaystyle \frac{p}{q}}+(1p)\mathrm{log}{\displaystyle \frac{1}{(1q)}}1`$
$``$ $`p\mathrm{log}{\displaystyle \frac{p}{q}}1`$
$`=`$ $`D(\rho |\sigma )1.`$
The first inequality follows from the Lindblad-Uhlmann monotonicity of relative entropy \[NC00\] and the second inequality follows because $`(1p)\mathrm{log}(1p)(\mathrm{log}e)/e>1`$, for $`0p1`$.
2. Define $`x_i=\mathrm{log}(p_i/q_i)`$. We can assume without loss of generality, by perturbing $`Q`$ slightly, that the values $`x_i`$ are distinct for distinct $`i`$. Let $`S^{}=\{i:x_i>0\}`$. Let $`D(P|Q)=k`$. Let
$$\text{positive real }l,S_l=\{i[n]:x_il\}.$$
Therefore,
$`k`$ $``$ $`\underset{P}{\mathrm{Pr}}[S_l]\mathrm{log}{\displaystyle \frac{\mathrm{Pr}_P[S_l]}{\mathrm{Pr}_Q[S_l]}}\underset{P}{\mathrm{Pr}}[S_l]l`$
$`\underset{P}{\mathrm{Pr}}[S_l]`$ $``$ $`k/l`$
Assume without loss of generality that $`x_1<x_2<\mathrm{}<x_n`$. Then if $`x_i>0`$, $`\mathrm{Pr}_P[S_{x_i}]k/x_i`$. Since $`S(P|Q)_{iS^{}}p_ix_i`$, the upper bound on $`S(P|Q)`$ is maximized when $`S^{}=\{2,\mathrm{},n\}`$, $`p_n=k/x_n`$, $`p_i=k(1/x_i1/x_{i+1})`$ for all $`i\{2,\mathrm{},n1\}`$, $`p_n=k/x_n`$, and $`p_1=1_{i=2}^np_i`$. Then,
$`S(P|Q)`$ $``$ $`{\displaystyle \underset{i=2}{\overset{n}{}}}p_ix_i`$
$`=`$ $`k{\displaystyle \underset{i=2}{\overset{n1}{}}}x_i(1/x_i1/x_{i+1})+k`$
$`=`$ $`k{\displaystyle \underset{i=2}{\overset{n1}{}}}{\displaystyle \frac{x_{i+1}x_i}{x_{i+1}}}+k`$
$``$ $`k{\displaystyle \underset{i=2}{\overset{n1}{}}}1+k`$
$`=`$ $`k(n1).`$
3. Let us measure $`\rho `$ and $`\sigma `$ in the eigenbasis of $`\sigma `$. We get two distributions: $`P_\rho `$ and $`P_\sigma `$. Let $`D(P_\rho |P_\sigma )=k`$. From Part 2, it follows,
$`k(n1)=D(P_\rho |P_\sigma )(n1)`$ $``$ $`S(P_\rho |P_\sigma )`$
$`=`$ $`\text{Tr }P_\rho \mathrm{log}P_\rho \text{Tr }P_\rho \mathrm{log}P_\sigma `$
$``$ $`\mathrm{log}n\text{Tr }P_\rho \mathrm{log}P_\sigma `$
$`=`$ $`\mathrm{log}n\text{Tr }\rho \mathrm{log}\sigma `$
$`=`$ $`\mathrm{log}n+S(\rho |\sigma )\text{Tr }\rho \mathrm{log}\rho `$
$``$ $`\mathrm{log}n+S(\rho |\sigma )`$
The second equality above holds since the measurement was in the eigenbasis of $`\sigma `$.
Thus
$$S(\rho |\sigma )D(P|Q)(n1)+\mathrm{log}n$$
4. Fix $`a>1`$, $`k>0`$. Let $`p_1=(a1)/a`$,
$$i\{2,\mathrm{},n1\}p_i=(a1)/a^i,$$
and $`p_n=1/a^{n1}`$. Let
$$i\{2,\mathrm{},n\}q_i=p_i/2^{ka^{i1}},$$
and $`q_1=1_{i=2}^nq_i`$. For any $`r>1`$, consider $`\stackrel{~}{P}=(p_1,\mathrm{},p_{\mathrm{log}_ar+1},0,\mathrm{},0)`$ normalized to make it a probability vector. It is easy to see that $`P\stackrel{~}{P}_12/r`$ and $`\frac{(r1)\stackrel{~}{P}}{r2^{rk}}Q`$. This shows that $`P,Q`$ satisfy the classical $`k`$-substate property, hence $`D(P|Q)2(k+1)`$.
Now,
$`S(P|Q)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{n}{}}}p_i\mathrm{log}{\displaystyle \frac{p_i}{q_i}}`$
$``$ $`p_1\mathrm{log}p_1+{\displaystyle \underset{i=2}{\overset{n}{}}}p_i\mathrm{log}{\displaystyle \frac{p_i}{q_i}}`$
$`>`$ $`1+(n2){\displaystyle \frac{k(a1)}{a}}+k`$
$`=`$ $`k(n1){\displaystyle \frac{k(n2)}{a}}1.`$
By choosing $`a`$ large enough, we can achieve $`S(P|Q)>k(n2)1`$. This shows that $`S(P|Q)>(D(P|Q)/21)(n2)1`$.
5. Proof skipped. Please see \[JRS02\] for a detailed proof.
6. Let $`M`$ be a POVM such that
$$k_1D(P|Q)=\text{Tr }MP\mathrm{log}\frac{\text{Tr }MP}{\text{Tr }MQ}$$
Let $`p\text{Tr }MP`$ and $`q\text{Tr }MQ`$. Therefore,
$$k_1=p\mathrm{log}\frac{p}{q}q=\frac{p}{2^{k_1/p}}$$
Let $`r=2/p`$. Since $`P`$ and $`Q`$ have the $`k`$-substate property, let $`P_r`$ be the distribution such that,
$$PP_r_t\frac{2}{r}=p\text{ (holds for classical distributions }\text{[JRS02]}\text{}$$
(1)
and
$$Q(1\frac{1}{r})\frac{P_r}{2^{rk}}0$$
(2)
Let $`p_r\text{Tr }MP_r`$. From (1) it follows,
$$p_r\frac{p}{2}$$
(3)
Also
$`{\displaystyle \frac{p}{2^{k_1/p}}}=q`$ $`=`$ $`\text{Tr }MQ`$
$``$ $`(1{\displaystyle \frac{1}{r}})\text{Tr }{\displaystyle \frac{MP_r}{2^{rk}}}\text{(from (}\text{2}\text{))}`$
$`=`$ $`(1{\displaystyle \frac{p}{2}}){\displaystyle \frac{p_r}{2^{rk}}}\text{(from definition)}`$
$``$ $`({\displaystyle \frac{1}{2}}){\displaystyle \frac{p_r}{2^{rk}}}\text{(since}p1\text{)}`$
$``$ $`({\displaystyle \frac{1}{2}}){\displaystyle \frac{p}{2^{rk+1}}}\text{(from (}\text{3}\text{))}`$
$`2^{rk+2}`$ $``$ $`2^{k_1/p}`$
$`rk+2`$ $``$ $`k_1/p`$
$`p(rk+2)`$ $``$ $`k_1`$
$`prk+2`$ $``$ $`k_1\text{(since}p1\text{)}`$
$`2k+2`$ $``$ $`k_1`$
Remark : This proof does not work for the quantum case because of $`\sqrt{(}r)`$ in the substate property.
7. $`S(\rho )`$ $`=`$ $`\text{Tr }\rho \mathrm{log}\rho \text{Tr }\rho \mathrm{log}\sigma `$
$``$ $`\text{Tr }\rho \mathrm{log}\rho \text{Tr }\rho \mathrm{log}{\displaystyle \frac{\rho }{2^k}}`$
$`=`$ $`k\text{Tr }\rho =k`$
The first inequality above follows from monotonicity of the operator log function.
8. We know that there exists a POVM element $`M`$ such that,
$$D(\rho |\sigma )=\text{Tr }M\rho \mathrm{log}\frac{\text{Tr }M\rho }{\text{Tr }M\sigma }$$
Let $`p=\text{Tr }M\rho `$ and $`q=\text{Tr }M\sigma `$. Let $`P=(p,1p)`$ and $`Q=(q,1q)`$. Note that
$`S(P|Q)`$ $`=`$ $`p\mathrm{log}p/q+(1p)\mathrm{log}(1p)/(1q)`$
$``$ $`p\mathrm{log}p/q1`$
$`=`$ $`D(\rho |\sigma )1`$
From Part 2 it follows that
$`S(\rho |\sigma )`$ $``$ $`D(\rho |\sigma )(n1)+\mathrm{log}n`$
$``$ $`(S(P|Q)+1)(n1)+\mathrm{log}n`$
$`S(P|Q)`$ $``$ $`{\displaystyle \frac{S(\rho |\sigma )\mathrm{log}n}{n1}}1`$
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# 1 Introduction
## 1 Introduction
Recently, there has been some interest in the question of whether or not defects (or impurities) of various kinds might be integrable and, if so, what kind of properties they might have. Already ten years ago it was pointed out by Delfino, Mussardo and Simonetti that the standard ideas of factorisation with a bulk, Lorentz invariant, S-matrix are incompatible (in most cases) with both reflection and transmission at a defect (the exceptions occurring when the S-matrix is $`\pm 1`$, or possibly when the defect has internal degrees of freedom ). Also, Konik and LeClair have examined the possibility of having a purely transmitting defect in the sine-Gordon model, in the sense of investigating algebraically the equations satisfied by the quantum transmission matrix. Some of their results appear to be relevant to the case of interest in this article although their starting point and emphasis were very different. More recently, some of the standard heuristic ideas have been questioned and modified within a different scheme developed by Mintchev, Ragoucy and Sorba and applied to various models . Other ideas can be found in and . However, the principal focus of those works has been to combine reflection and transmission, whereas the present article will focus exclusively on purely transmitting integrable defects of a quite particular type.
In it was noted, in a Lagrangian approach, that a field theory could permit a discontinuity or ‘jump-defect’ and yet remain classically integrable. The principal examples of this phenomenon discovered so far have been very specific. These include free scalar fields, Liouville theory, the sine-Gordon model, and certain affine Toda models , all of which permit Bäcklund transformations. Indeed, the defect conditions relating fields evaluated as limits from either side of the jump-defect turn out to be a Bäcklund transformation frozen at the location of the defect. This fact is quite striking in view of the importance Bäcklund transformations have played in the development of soliton theory. By now there is an extensive literature on Bäcklund transformations, and their uses, in a variety of contexts; see for example . A particular feature of this kind of jump-defect is precisely that it is purely transmitting at the classical level, and presumably also at the quantum level. In a sense, a jump-defect is more severe than the more usual model of an impurity represented by a delta-function contribution in the equations of motion, since the latter generally requires the field, though not necessarily its spatial derivatives, to be continuous. On the other hand, a defect of delta-function type does not generally maintain classical integrability, although there are a variety of other interesting phenomena involving solitons associated with it . A jump-defect is also simpler because, in the examples found so far, the classical systems have no periodic solutions specifically associated with the presence of the defect. This is in contrast to integrable field theory with a boundary in which boundary breathers may be found for a suitable range of parameters even in the simplest of models .
It is the purpose of this article to investigate the quantum sine-Gordon field theory in the presence of a jump-defect and to locate within the quantum framework the influence of at least some of the striking consequences of the classical model. Details are provided in but for convenience, a brief summary will also be provided in the next section. Particular attention will be paid to the properties of solitons scattering with the defect. For reviews of models with solitons in general, including the sine-Gordon model, see for example ; much essential information concerning the quantum sine-Gordon model is to be found in the review by Zamolodchikov and Zamolodchikov .
## 2 Classical sine-Gordon with a jump-defect
The sine-Gordon model in the bulk will be taken to be defined by the Lagrangian density
$$=\frac{1}{2}\left((_tu)^2(_xu)^2\right)\frac{m^2}{\beta ^2}(1\mathrm{cos}\beta u),$$
(2.1)
although for classical considerations it is often convenient to remove the mass parameter $`m`$ and the coupling $`\beta `$ by a rescaling. A single jump-defect placed at $`x=0`$ is described by modifying the Lagrangian in the following manner, where the field on the left of the defect is denoted by $`u`$ and the field on the right of it by $`v`$. The full Lagrangian will consist of pieces from the bulk regions ($`x<0`$ and $`x>0`$), together with a delta function contribution at $`x=0`$. In detail, the Lagrangian density is given by
$$=\theta (x)_u+\theta (x)_v\delta (x)\left[\frac{1}{2}\left(uv_tvu_t\right)(u,v)\right]$$
(2.2)
with
$$=\frac{2m\sigma }{\beta ^2}\mathrm{cos}\beta \left(\frac{u+v}{2}\right)\frac{2m}{\sigma \beta ^2}\mathrm{cos}\beta \left(\frac{uv}{2}\right).$$
(2.3)
Following from this the (suitably scaled) bulk field equations and defect conditions are:
$`x<0:^2u=\mathrm{sin}u,`$
$`x>0:^2v=\mathrm{sin}v,`$
$`x=0:_xu_tv=\sigma \mathrm{sin}\left({\displaystyle \frac{u+v}{2}}\right){\displaystyle \frac{1}{\sigma }}\mathrm{sin}\left({\displaystyle \frac{uv}{2}}\right)`$
$`_xv_tu=\sigma \mathrm{sin}\left({\displaystyle \frac{u+v}{2}}\right){\displaystyle \frac{1}{\sigma }}\mathrm{sin}\left({\displaystyle \frac{uv}{2}}\right).`$ (2.4)
The term containing first order time derivatives in the defect Lagrangian is required by integrability and has interesting consequences. Clearly, it is not invariant under reversing the sense of time, nor is it invariant under the bulk symmetries $`uu+2a\pi /\beta ,vv+2b\pi /\beta `$, where $`a`$ and $`b`$ are integers; it is not even invariant under the reflections $`uu`$ or $`vv`$. However, it is invariant under certain combinations, such as reflecting both fields simultaneously or reflecting one of them and reversing the sense of time. Under the bulk symmetries, which translate among the different bulk ground states, the defect term changes by a total time derivative (note that the defect potential itself is invariant under these transformations only if $`a\pm b`$ is even). However, that is insufficient for invariance under all circumstances because how the action changes will depend upon boundary conditions - in this sense it is reminiscent of a one-dimensional ‘Chern-Simons’ term - and these may vary in the presence of scattering configurations of fields. Additional comments concerning this aspect in the quantum field theory will be made at the end of section (3).
The Lagrangian does not violate time translation and, provided a contribution from the defect, involving the fields evaluated at the defect, is included, the total energy,
$$=E(u)+E(v)+,$$
will be conserved. Moreover, setting
$$\frac{𝒰}{u}=\frac{}{v},\frac{𝒰}{v}=\frac{}{u},$$
so that,
$$𝒰=2\sigma \mathrm{cos}\left(\frac{u+v}{2}\right)+\frac{2}{\sigma }\mathrm{cos}\left(\frac{uv}{2}\right),$$
one finds that the total momentum,
$$𝒫=P(u)+P(v)+𝒰,$$
is also conserved . Hence the fields on either side of $`x=0`$ can exchange both energy and momentum with the defect. In fact, the conditions (2.2) are immediately recognisable as a Bäcklund transformation frozen at $`x=0`$ and it would be attractive to have a specific physical mechanism to implement it. It is worth emphasising that the defect potential has twice the period of the bulk potentials, a feature it shares with the potentials representing integrable boundary conditions , which restrict the sine-Gordon model to a half-line without destroying its integrability. It has already been noted that the defect Lagrangian is not invariant under time-reversal or parity. Rather there is in each case an extended symmetry involving the fields and the parameter $`\sigma `$. In the case of parity this will be further remarked upon below.
In the sine-Gordon model a single soliton may be described conveniently by
$$e^{iu/2}=\frac{1+iE}{1iE},E=e^{ax+bt+c},$$
(2.5)
where $`a,b,e^c`$ are all real, with $`a^2b^2=1`$, and it is useful to parameterize $`a`$ and $`b`$ in terms of rapidity, setting $`a=\mathrm{cosh}\theta ,b=\mathrm{sinh}\theta `$. With $`\theta >0`$, the soliton is moving along the $`x`$-axis in a positive direction. An anti-soliton (having the same velocity and location) is obtained by the replacement $`EE`$ (or equivalently $`cc+i\pi `$). The mass of the soliton is 8 in the units for which the fields $`u`$ or $`v`$ have unit mass parameter. \[If the coupling and mass scale are reinserted, the mass parameter of the fields will be $`m`$ and that of the soliton will be $`8m/\beta ^2`$.\]
Supposing a soliton given by (2.5), moving in a positive sense along the $`x`$-axis, encounters the defect, then a similar, but delayed, soliton emerges given by
$$e^{iv/2}=\frac{1+izE}{1izE},E=e^{ax+bt+c},$$
(2.6)
where $`z`$ represents the delay . Using the defect conditions,
$$z=\frac{e^\theta +\sigma }{e^\theta \sigma },$$
(2.7)
and there will be a variety of possible consequences according to the choice of $`\sigma `$. Setting $`\sigma =e^\eta `$, the expression for $`z`$ can be written alternatively as
$$z=\mathrm{coth}\left(\frac{\eta \theta }{2}\right).$$
If $`\theta >0`$ and if $`\eta <0`$ it is clear $`z<0`$, implying that the incoming soliton always converts to an anti-soliton, although for large $`\theta `$ it will be delayed very little. On the other hand, for $`\eta >0`$ there are several possibilities: if $`\theta <\eta `$ the soliton will be delayed but its character remains unchanged; if $`\theta >\eta `$ the soliton flips to an anti-soliton; but if $`\theta =\eta `$ the incoming soliton is infinitely delayed and therefore never emerges from the defect. In effect, a soliton, which in the distant past interpolated between $`0`$ and $`2\pi `$, is replaced in the far future by the static solution $`u=0,v=2\pi `$. It is easy to check that the latter satisfies the defect conditions by itself and stores at the defect location precisely the energy and momentum originally carried by the soliton. A soliton travelling in the opposite direction will be affected in a similar variety of ways. Several solitons passing the same defect are each delayed by a similar factor. In effect, the multi-soliton solution $`u`$ constructed by assembling, using Hirota’s method , a set of exponential factors
$$E_k=e^{a_kx+b_kt+c_k},$$
is replaced by a similar solution $`v`$ in which each exponential factor is multiplied by the appropriate factor $`z_k`$, given by (2.7) with $`\theta `$ replaced by $`\theta _k`$.
There might be many defects placed at different locations, each with its own parameter, and each affecting a passing soliton independently of all the others. In particular, a soliton passing two defects, which are separated spatially but otherwise identical, always retains its character (since if it is flipped by one it must be flipped by the other). However, it will be delayed by the combined factor
$$z=\mathrm{coth}^2\left(\frac{\eta \theta }{2}\right).$$
Interestingly, this is precisely the delay it would have experienced had it been overtaken (or been overtaken by) a soliton of rapidity $`\eta `$ . Moreover, a current of such solitons would build up positive and negative topological charges on the two defects - indicating that a pair of similar defects might behave like a store of topological charge. In other words, a pair of similar defects might be thought of as an analogue of a capacitor in an electrical circuit.
The behaviour of solitons in the presence of jump-defects suggests that if mathematical defects such as these can be found within actual physical systems then using them to control solitons might lead to technological applications, perhaps along the lines suggested in . They also suggest that the jump-defect behaves, in a sense, as though it were ‘half’ a soliton. Even the energy or momentum naturally associated with a defect via the expressions for the total energy and momentum relate to an object of mass $`4m/\beta ^2`$ rather than one of mass $`8m/\beta ^2`$.
Since any number of jump-defects can co-exist, each influencing the progress of a soliton independently of all the others, and since they are all at rest, it is clear the defects do not exert any nett long-range influence on each other. In this respect, the defects are quite different to classical sine-Gordon solitons since multi-soliton solutions with each soliton at rest at an arbitrarily chosen location cannot exist. Famously, there are extended structures in three spatial dimensions which may be assembled to create multi-particle-like stationary solutions to their equations of motion. Yang-Mills-Higgs BPS-monopoles provide a prime example of this phenomenon (for a recent review, see ): classical solutions can be constructed with arbitrary numbers of like-charged magnetic monopoles balancing at rest despite the existence of long-range forces between them. In this case, there are two kinds of long range force (the electromagnetic force and the Higgs force), exactly cancelling out. There is no reason, in principle, why the defects should not move, and some comments concerning that possibility will be made in section(6). If they are able to move with different speeds then inevitably they must scatter and the most interesting question concerns the nature of this scattering and the nature of the associated short-range interaction. A further question will be whether the defects themselves are describable by a quantum field theory.
For future reference, it is also instructive to calculate the transmission factor through the defect in the situation where the equations (2.2) are linearized. If the linear perturbation is a perturbation of the static solution $`u=2n\pi ,v=2m\pi `$, then it will be denoted $`T_{\mathrm{even}}`$ or $`T_{\mathrm{odd}}`$ according to whether $`nm`$ is even or odd. It is straightforward in either case to show there is no reflection and $`u`$ and $`v`$ will have the form
$$u=e^{i\omega t+ikx},x<0;v=Te^{i\omega t+ikx},x>0$$
(2.8)
with
$$T_{\mathrm{even}}(\theta ,\eta )=i\frac{\mathrm{sinh}\left(\frac{\theta \eta }{2}\frac{i\pi }{4}\right)}{\mathrm{sinh}\left(\frac{\theta \eta }{2}+\frac{i\pi }{4}\right)}=\overline{T}_{\mathrm{odd}}(\theta ,\eta ).$$
(2.9)
In the limit $`\eta \mathrm{}`$, (or $`\sigma 0`$), the transmission factor tends to unity, as it should since continuity of the fields is restored and the jump disappears. Notice that the transmission factors are unitary but they do not satisfy $`T(\theta )T(\theta )=1`$. This is a consequence of the behaviour of the defect conditions under parity or time reversal, neither of which is an invariance of (2.2). For example, a parity transformation interchanges $`u`$ and $`v`$ but needs to be accompanied by the change of sign of one of the fields $`u`$ or $`v`$ and a replacement of $`\sigma `$ by $`1/\sigma `$, or equivalently, $`\eta \eta `$. Thus, defining $`T_P(\theta ,\eta )=T(\theta ,\eta )`$ one finds
$$T(\theta )T_P(\theta )=1.$$
(2.10)
However, for a specific $`\sigma `$ both parity and time-reversal are explicitly broken.
In the quantum field theory it is expected that the bound states, or ‘breathers’ will be transmitted through a defect and suffer a change of phase whose classical limit should be one of $`T_{\mathrm{even}}`$ or $`T_{\mathrm{odd}}`$, depending on the precise circumstances.
A couple of other observations are in order. It has already been pointed out that the defect conditions (2.2) have the form of a Bäcklund transformation with the spatial derivatives fixed at the location of the defect. In the bulk, a Bäcklund transformation generates solitons in the sense that if one of the fields is taken to be zero and the equation for the other is integrated, the result is typically a one-soliton solution. Similarly, if the first is taken to be a one-soliton, integrating for the second will give a two-soliton, and so on . The question is: how does this behaviour fit in with the defect conditions given above? Clearly, since the defect conditions are a frozen Bäcklund transformation, there is no question of integrating the equations for one field given the other. Up to this point the conditions have been used to demonstrate how a single soliton approaching the defect will be delayed on passing through, or possibly altered more drastically. However, the possibility of two solitons emerging has not been considered.
Consider the possibility that (2.2) allows a single soliton described by $`u`$ to approach the defect, and a two-soliton, described by $`v`$ to emerge from it. It is not difficult to check that for this to happen one of the emerging solitons must be a delayed version of the original one and the new soliton has parameters related to the defect parameter $`\sigma `$ (although there is no information that would fix its ‘position’). If this is possible then as $`t\mathrm{}`$, $`E(v)`$ must have a contribution from both solitons, whereas as $`t\mathrm{}`$, $`E(u)`$ has a contribution merely from one of them. To balance the energy, the additional contribution must come from the difference of the energies initially and finally stored in the defect. Supposing $`\sigma >0`$, and initially there is no discontinuity at the defect (that is, $`u(0,\mathrm{})=2\pi =v(0,\mathrm{}))`$, then the defect contribution to the energy (2.3) will be negative and equal to
$$2\left(\sigma +\frac{1}{\sigma }\right)=4\mathrm{cosh}\eta .$$
(2.11)
On the other hand, if subsequently there are two solitons on the right, the defect must end up with a $`2\pi `$ discontinuity, meaning the energy stored there must have increased since it will ultimately have to be
$$2\left(\sigma +\frac{1}{\sigma }\right)=4\mathrm{cosh}\eta .$$
(2.12)
Since the additional soliton also contributes positively to the energy the overall energy conservation is violated. So, this situation cannot actually occur. The only possibility in these circumstances is for the approaching soliton to pass through - albeit with a delay (or conversion to an anti-soliton). Hence, no energy is extracted from the defect and the final configuration of the fields at the defect is $`u(0,\mathrm{})=0=v(0,\mathrm{})`$. Another possible starting configuration has a $`2\pi `$ discontinuity at the defect, meaning a store of positive energy which is exactly right to allow a new soliton (or anti-soliton) to emerge in the final state, leaving a $`4\pi `$ (or $`0`$) discontinuity behind. The new soliton will have rapidity $`\theta =\eta `$ and energy $`E=8\mathrm{cosh}\eta `$. However, the classical system provides no information as to the relative position, or character, of the additional soliton.
Since $`\sigma `$ is a free parameter, it is also possible to take $`\sigma <0`$. Then the energy stored in the defect would be positive to begin with provided the initial discontinuity was zero modulo $`2\pi `$. Then, in principle, a soliton could emerge making use of that energy leaving a discontinuity behind.
Thus it appears defects might produce solitons as well as absorbing them. However, since there is no information concerning the location (or time-origin) of the additional soliton the situation bears a resemblance to an excited atom. Thought of classically, there is no information to indicate the time of decay. Instead, quantum mechanics is needed to supply a probability of decay within a specified time. This analogy suggests the quantum story of defects of the kind considered here could be considerably more interesting and motivates a search for the quantum analogues of the transmission matrices.
To date, the discussion has been entirely theoretical and it would be even more interesting to find a physical situation in which Bäcklund transformations play a substantive role, rather than being a mathematical, solution-generating, device. In such a situation it would be expected that the effects described briefly above should be manifest and amenable to observation.
## 3 Transmission matrices
Rather than depending upon the previous literature for the transmission matrices they will be derived afresh from first principles, guided by the classical features of the jump-defect. In any case, there will be some significant differences with respect to earlier work on this topic.
On the basis of the classical scattering of a soliton from the jump-defect it is expected the defect will be purely transmitting and able to store topological charge. Moreover, topological charge may be added or removed in steps of two. The classical picture suggests there will be two types of transmission matrix depending on whether the defect is carrying an even or odd charge. If $`\sigma >0`$, only the even transmission matrix is expected to be unitary since the ‘vacuum’, or least energy configuration of a defect with a positive parameter $`\sigma `$, must have even (positive, negative, or zero) topological charge and cannot decay. Thus, there should be a transmission matrix
$${}_{}{}^{\mathrm{e}}T_{a\alpha }^{b\beta }(\theta ),$$
regarded as describing the transmission from the region $`x<0`$ to the region $`x>0`$, where $`a`$ and $`b`$ may be $`+`$ for a soliton or $``$ for an anti-soliton, and the labels $`\alpha `$ and $`\beta `$ are even integers (positive, negative or zero). These transmission matrices should satisfy (for real rapidity $`\theta `$),
$$^\mathrm{e}T(\theta )^\mathrm{e}T^{}(\theta )=1.$$
(3.1)
More precisely, the transmission matrices relate states of the system in the far future to those in the far past and the state of the system is labelled by the soliton energy-momentum, or rapidity, and its character (topological charge) together with the state of the defect labelled by the topological charge accumulated on it. On the other hand, because of the properties of the Lagrangian under a parity transformation the transmission matrix for a soliton moving in the opposite direction will be different.
The transmission matrix will depend also on the defect parameter and the bulk coupling as well as rapidity, and satisfy a number of other relations to be detailed below.
The soliton-soliton bulk scattering matrix is taken to be the standard one , given by
$$S_{kl}^{mn}(\mathrm{\Theta })=\rho (\mathrm{\Theta })\left(\begin{array}{cccc}a(\mathrm{\Theta })& 0& 0& 0\\ 0& c(\mathrm{\Theta })& b(\mathrm{\Theta })& 0\\ 0& b(\mathrm{\Theta })& c(\mathrm{\Theta })& 0\\ 0& 0& 0& a(\mathrm{\Theta })\end{array}\right),$$
(3.2)
where $`k,l`$ label the incoming particles and $`m,n`$ label the outgoing particles in a two-body scattering process, with the particles labelled $`k,n`$ having rapidity $`\theta _1`$, and the particles labelled $`l,m`$ having rapidity $`\theta _2`$. The various pieces of the matrix are defined by
$$a(\mathrm{\Theta })=\frac{qx_2}{x_1}\frac{x_1}{qx_2},b(\mathrm{\Theta })=\frac{x_1}{x_2}\frac{x_2}{x_1},c(\mathrm{\Theta })=q\frac{1}{q},$$
(3.3)
with
$$\mathrm{\Theta }=\theta _1\theta _2,q=e^{i\pi \gamma },x_p=e^{\gamma \theta _p}.$$
(3.4)
In this notation the crossing property of the S-matrix is represented by
$$S_{kl}^{mn}(i\pi \mathrm{\Theta })=S_{km}^{ln}(\mathrm{\Theta }),$$
(3.5)
with the diagonal elements $`S_+^+(\mathrm{\Theta })`$ and $`S_+^+(\mathrm{\Theta })`$ crossing into themselves. The overall factor $`\rho (\mathrm{\Theta })`$ will be needed later and is given by:
$$\rho (\mathrm{\Theta })=\frac{\mathrm{\Gamma }(1+i\gamma \mathrm{\Theta }/\pi )\mathrm{\Gamma }(1\gamma i\gamma \mathrm{\Theta }/\pi )}{2\pi i}\underset{k=1}{\overset{\mathrm{}}{}}R_k(\mathrm{\Theta })R_k(i\pi \mathrm{\Theta }),$$
(3.6)
where
$$R_k(\mathrm{\Theta })=\frac{\mathrm{\Gamma }(2k\gamma +i\gamma \mathrm{\Theta }/\pi )\mathrm{\Gamma }(1+2k\gamma +i\gamma \mathrm{\Theta }/\pi )}{\mathrm{\Gamma }((2k+1)\gamma +i\gamma \mathrm{\Theta }/\pi )\mathrm{\Gamma }(1+(2k1)\gamma +i\gamma \mathrm{\Theta }/\pi )}.$$
(3.7)
Note, the conventions adopted by Konik and LeClair have been used. Therefore, in particular, the coupling $`\gamma `$ in terms of the Lagrangian coupling $`\beta `$ (with $`\mathrm{}=1`$ and the conventions indicated by (2.1)) is defined by
$$\frac{1}{\gamma }=\frac{\beta ^2}{8\pi \beta ^2}.$$
(3.8)
Since the defect is purely transmitting, the usual heuristic arguments based on factorisability and bulk integrability would require
$$S_{kl}^{mn}(\mathrm{\Theta })^\mathrm{e}T_{n\alpha }^{t\beta }(\theta _1)^\mathrm{e}T_{m\beta }^{s\gamma }(\theta _2)=^\mathrm{e}T_{l\alpha }^{n\beta }(\theta _2)^\mathrm{e}T_{k\beta }^{m\gamma }(\theta _1)S_{mn}^{st}(\mathrm{\Theta }),$$
(3.9)
and this seems to be the most appropriate assumption to make in the present context. There is also a transmission matrix representing transmission through the defect of a soliton moving from right to left. However, this will be determined in terms of $`{}_{}{}^{\mathrm{e}}T`$ and part of the purpose in this section is to develop a set of criteria constraining $`{}_{}{}^{\mathrm{e}}T`$ itself without reference to transmission from right to left. This means that the important crossing properties of $`{}_{}{}^{\mathrm{e}}T`$ will not be part of the initial story. This may appear to be an unconventional way to proceed but it is useful to disentangle those aspects of the scheme which are purely algebraic, or depend upon conjectured general principles, from those which depend upon the special nature of the Lagrangian model defined via (2.1) and (2.2).
The use of a bulk scattering matrix depending only on the rapidity difference might be questioned in a situation which appears manifestly to break Lorentz invariance. However, to decide that question will require a full discussion of moving defects going beyond the scope of the present paper. For the time being (3.9) will be used as a working hypothesis to be abandoned, if necessary, when the details emerge of a much fuller picture. This will include the scattering of defects themselves. One indication that the defect will itself behave like a particle lies in the results obtained below where the transmission matrices definitely depend upon the difference of $`\theta `$, the rapidity of a soliton, and $`\eta `$, which in view of the classical defect conditions being a frozen Bäcklund transformation has the character of a rapidity. Indeed, a Lorentz transformation on the system will be compensated by a change of the parameter $`\sigma `$ and a shifting in the location of the defect. The scattering of defects is likely to be entirely consistent classically owing to Bianchi’s celebrated theorem of permutability for Bäcklund transformations . Further remarks on this will be made in section (6) where a start on the construction of the quantum scattering matrix will also be given.
Equations (3.9) may have several solutions and the purpose of this article is to find among these one which matches qualitatively the jump-defect situation and to find ways to accumulate evidence for it. For some purposes it is convenient to change notation slightly and write
$$^\mathrm{e}T=\left(\begin{array}{cc}T_+^+& T_+^{}\\ T_{}^+& T_{}^{}\end{array}\right)\left(\begin{array}{cc}A& B\\ C& D\end{array}\right),$$
(3.10)
where the block matrix entries are labelled by the topological charge of the defect (and are therefore infinite dimensional).
For a first examination, it is useful to make use of the topological charge conservation to note that $`A,D`$ are diagonal while $`B,C`$ are slightly off-diagonal:
$$A_\alpha ^\beta =a_\alpha \delta _\alpha ^\beta ,D_\alpha ^\beta =d_\alpha \delta _\alpha ^\beta ,B_\alpha ^\beta =b_\alpha \delta _\alpha ^{\beta 2},C_\alpha ^\beta =c_\alpha \delta _\alpha ^{\beta +2}.$$
(3.11)
Then, the transmission relations reduce to a number of equations among the matrices $`A,B,C,D`$ and the entries of the $`S`$-matrix, which fall into three types. The first is a set of four two-term relations (for the purposes of these the subscripts or superscripts ‘1’ and ‘2’ refer to the rapidities of the incoming particles):
$$A_1A_2=A_2A_1,D_1D_2=D_2D_1,B_1B_2=B_2B_1,C_1C_2=C_2C_1.$$
(3.12)
The first two are automatically satisfied since $`A`$ and $`D`$ are diagonal but the other pair provide genuine constraints. Consider the third: in terms of components one finds
$$b_\alpha ^1b_{\alpha +2}^2=b_\alpha ^2b_{\alpha +2}^1\text{or}\frac{b_{\alpha +2}^1}{b_\alpha ^1}=\frac{b_{\alpha +2}^2}{b_\alpha ^2},$$
(3.13)
and the second of these implies that neither side can depend on the rapidity, implying
$$b_\alpha (\theta )=\sigma ^\alpha b_0(\theta ),$$
(3.14)
where $`\sigma `$ has no dependence on rapidity. In fact, there are two such solutions according to whether $`\alpha `$ is selected to be even or odd. Similarly, the solution for the components of $`C`$ are:
$$c_\alpha (\theta )=\tau ^\alpha c_0(\theta ),$$
(3.15)
where $`\tau `$ has no dependence on rapidity. A second group of four equations has the form
$`b(B_1C_2C_2B_1)`$ $`=`$ $`c(D_2A_1D_1A_2)`$
$`b(B_2C_1C_1B_2)`$ $`=`$ $`c(A_1D_2A_2D_1)`$
$`b(A_1D_2D_2A_1)`$ $`=`$ $`c(C_2B_1C_1B_2)`$
$`b(A_2D_1D_1A_2)`$ $`=`$ $`c(B_1C_2B_2C_1).`$ (3.16)
Since $`A,D`$ are diagonal, the last pair relate $`B`$ and $`C`$:
$$\frac{c_{\alpha +2}^1}{b_\alpha ^1}=\frac{c_{\alpha +2}^2}{b_\alpha ^2},$$
(3.17)
implying that neither side depends on $`\theta `$. Hence, to ensure this works for any $`\alpha `$ requires a constraint. Introducing a $`\theta `$independent constant $`\mu `$ this is
$$c_0(\theta )=\mu b_0(\theta ).$$
(3.18)
Using (3.17) the difference of the first two of (3) becomes an identity leaving one equation to investigate later. The third set of eight equations have the form,
$`aA_1B_2=bB_2A_1+cA_2B_1,`$ $`aB_1A_2=bA_2B_1+cB_2A_1,`$
$`aA_2C_1=bC_1A_2+cA_1C_2,`$ $`aC_2A_1=bA_1C_2+cC_1A_2,`$
$`aD_2B_1=bB_1D_2+cD_1B_2,`$ $`aB_2D_1=bD_1B_2+cB_1D_2,`$
$`aD_1C_2=bC_2D_1+cD_2C_1,`$ $`aC_1D_2=bD_2C_1+cC_2D_1.`$ (3.19)
Using (3.11), the first pair of these can be combined to give
$$\left(a\frac{a_\alpha ^1}{a_{\alpha +2}^1}b\right)\left(a\frac{a_{\alpha +2}^2}{a_\alpha ^2}b\right)=c^2,$$
(3.20)
which, on using the definitions of $`a,b,c`$, implies
$$\frac{a_{\alpha +2}}{a_\alpha }=q\text{or}1/q.$$
(3.21)
For either of the two choices there will be a solution for $`A`$. For example, choosing the first and setting $`q=Q^2`$ the solution has the form,
$$a_\alpha (\theta )=Q^\alpha a_0(\theta ).$$
(3.22)
Actually, since $`\alpha `$ may be a positive or negative integer both possibilities given in (3.21) are covered by (3.22). Inserting this in the first of the first pair in (3) reveals
$$x_2\frac{a_0^2}{b_0^2}=x_1\frac{a_0^1}{b_0^1}$$
(3.23)
implying that
$$b_0(\theta )=\lambda x_\theta a_0(\theta ),$$
(3.24)
with $`\lambda `$ independent of rapidity. Hence also (slightly redefining $`\mu `$),
$$c_0(\theta )=\mu x_\theta a_0(\theta ).$$
(3.25)
The second pair of equations involving $`A,C`$ is now an identity. The two equations of the third pair when combined reveal a similar relation to (3.20) and therefore $`D`$ can also be written in the form
$$d_\alpha (\theta )=Q^{ϵ\alpha }d_0(\theta ),ϵ=\pm 1.$$
(3.26)
In addition, the first of the third pair also reveals,
$$d_0(\theta )=\nu x_\theta ^{1+ϵ}a_0(\theta ),ϵ=\pm 1.$$
(3.27)
With these expressions for $`A,B,C,D`$, each of (3) is an identity. Finally, it is necessary to return to the first of the four-term relations (3). It now implies
$$bx_1x_2(\sigma \tau )^\alpha \mu \lambda \left(\tau ^2\frac{1}{\sigma ^2}\right)=cQ^{(1+ϵ)\alpha }\nu \left(x_2^{1+ϵ}x_1^{1+ϵ}\right).$$
(3.28)
If $`ϵ=1`$ then $`\tau ^2=1/\sigma ^2`$ is required but there are no other constraints on the coefficients. On the other hand, if $`ϵ=1`$ one requires,
$$\sigma \tau =Q^2,\nu =\frac{q\lambda \mu }{\sigma ^2}.$$
(3.29)
In view of (3.29) it would be convenient in this case to put $`\sigma =\rho Q,\tau =Q/\rho `$, with $`\rho `$ a free parameter, and then
$$\nu =\frac{\lambda \mu }{\rho ^2}.$$
(3.30)
At this point it is convenient to summarize the possibilities as follows:
$`a_\alpha =Q^\alpha A(\theta ),d_\alpha =\nu Q^{ϵ\alpha }x_\theta ^{1+ϵ}A(\theta ),ϵ=\pm 1`$
$`b_\alpha =\lambda \sigma ^\alpha x_\theta A(\theta ),c_\alpha =\mu \sigma ^\alpha x_\theta A(\theta ),\text{if}ϵ=1`$
$`b_\alpha =\lambda (\rho Q)^\alpha x_\theta A(\theta ),c_\alpha =\mu (Q/\rho )^\alpha x_\theta A(\theta ),\nu =\mu \lambda /\rho ^2\text{if}ϵ=+1.`$ (3.31)
The quantities $`\mu ,\nu ,\lambda `$ and $`A(\theta )`$ will differ between the two alternative choices of $`ϵ`$. Notice that the whole process, starting from (3.9) could be repeated assuming the defect labels on the $`T`$-matrix were odd since it transpires the even and odd solutions never mix via (3.9). On the other hand, for the reasons mentioned earlier, the transmission matrix with odd labels $`{}_{}{}^{\mathrm{o}}T(\theta )`$ refers to a situation which is unstable and, even if it can be defined, its properties are likely to differ from those of $`{}_{}{}^{\mathrm{e}}T(\theta )`$. The expressions (3) are similar to, but not quite the same as, the ansatz proposed by Konik and LeClair .
The usual bulk bound-state bootstrap operates, presumably, just with the left to right transmission matrix and will offer some further constraints . Suppose the pair of particles $`a,b`$ can form a bound state $`c`$ at the rapidities (usual conventions)
$$\theta _a=\theta _c+i\overline{U}_{a\overline{c}}^{\overline{b}},\theta _b=\theta _ci\overline{U}_{b\overline{c}}^{\overline{a}},\overline{U}=\pi U,$$
(3.32)
then the heuristics will demand that the transmission matrices are consistent with this. Thus,
$$c_{ab}^f{}_{}{}^{\mathrm{e}}T_{f\alpha }^{c\beta }(\theta _c)=^\mathrm{e}T_{b\alpha }^{d\gamma }(\theta _b)^\mathrm{e}T_{a\gamma }^{e\beta }(\theta _a)c_{de}^c,$$
(3.33)
where the coupling constants are denoted $`c_{ab}^f`$, and repeated indices are summed (as usual). Typically, this can be used to generate the transmission matrices for breathers and indeed that calculation will be performed later. However, it can also be used to discuss the ‘annihilation pole’ at which a particle and anti-particle virtually annihilate to the vacuum. This happens when $`b=\overline{a}`$ and $`\theta _a=\theta +i\pi /2,\theta _{\overline{a}}=\theta i\pi /2`$, and requires,
$$\delta _\alpha ^\beta =\underset{c}{\overset{\mathrm{e}}{}}T_{\overline{a}\alpha }^{c\gamma }\left(\theta \frac{i\pi }{2}\right)^\mathrm{e}T_{a\gamma }^{\overline{c}\beta }\left(\theta +\frac{i\pi }{2}\right).$$
(3.34)
Since the only possibility of annihilating to the vacuum in the present case must involve a soliton and an anti-soliton, the couplings $`c_{a\overline{a}}^0,c_{c\overline{c}}^0`$ cancel. To see the consequences of (3.34) it is convenient to shift $`\theta `$ by $`i\pi /2`$ and note $`x(\theta +i\pi )=qx(\theta )`$. Then, with the components of the transmission matrix given by the expressions (3.11) and (3), the off-diagonal components of (3.34) are identically satisfied for either choice of $`ϵ`$, while the two diagonal entries cannot be satisfied for the choice $`ϵ=1`$. For the other choice, $`ϵ=1`$, the diagonal entries each lead to the same condition, namely,
$$\left(\nu +\frac{\lambda \mu qx^2}{\sigma ^2}\right)A(\theta )A(\theta +i\pi )=1.$$
(3.35)
Hence, setting $`A(\theta )=f(q,x)/\sqrt{\nu }`$, one needs to solve
$$(1+p^2x^2)f(q,x)f(q,qx)=1,p^2=\lambda \mu q/\nu \sigma ^2.$$
(3.36)
At this point it is worth returning to the unitarity condition, equation (3.1), which, on choosing the above solutions with $`ϵ=1`$, requires
$$\overline{\sigma }=1/\sigma ,\lambda =\overline{\mu }\nu \sigma ^2/q,p^2=\overline{\mu }\mu ,\overline{\nu }\nu =1,\overline{A}(\theta )A(\theta )(1+\overline{\mu }\mu x^2)=1.$$
(3.37)
This follows immediately and the details will be omitted. Thus, in terms of $`f(q,x)`$, the two equations which must be satisfied simultaneously are:
$$f(q,x)f(q,qx)(1+p^2x^2)=1=\overline{f}(q,x)f(q,x)(1+p^2x^2),pxe^{\gamma (\theta \eta )}e^{\gamma \stackrel{~}{\theta }}.$$
(3.38)
Writing $`p`$ in terms of $`\eta `$ is deliberate and anticipates the eventual identification of this parameter with the defect parameter described in the introduction. As an aside, for $`ϵ=1`$ the off-diagonal terms in (3.1) cannot be satisfied, and therefore this possibility is ruled out by the unitarity condition alone. Henceforth, for both of the above reasons, it will be assumed $`ϵ=1`$.
Clearly, (3.38) implies
$$\overline{f}(q,x)=f(q,qx),$$
(3.39)
and in essence these are the same equations as those solved by Konik and LeClair. It is convenient to write
$$f(q,x)=\frac{e^{\gamma \stackrel{~}{\theta }/2}e^{i\pi \gamma /4}}{\sqrt{2\pi }}g(q,x)$$
(3.40)
so that
$$g(q,x)\overline{g}(q,x)=\frac{\pi }{\mathrm{cosh}\gamma \stackrel{~}{\theta }},\overline{g}(q,x)=g(q,qx).$$
(3.41)
A solution to the latter pair of equations is given by
$$g(q,x)=\mathrm{\Gamma }(1/2i\gamma \stackrel{~}{\theta }/\pi )\underset{k=1}{\overset{\mathrm{}}{}}\frac{\mathrm{\Gamma }(1/2+2k\gamma i\gamma \stackrel{~}{\theta }/\pi )\mathrm{\Gamma }(1/2+(2k1)\gamma +i\gamma \stackrel{~}{\theta }/\pi )}{\mathrm{\Gamma }(1/2+2k\gamma +i\gamma \stackrel{~}{\theta }/\pi )\mathrm{\Gamma }(1/2+(2k1)\gamma i\gamma \stackrel{~}{\theta }/\pi )}.$$
(3.42)
On the other hand, the ‘minimal’ solution given by Konik and LeClair can be written in the following way:
$$f(q,x)=\frac{e^{i\pi (1+\gamma )/4}}{(1+ipx)}\frac{r(x)}{\overline{r}(x)},$$
(3.43)
where
$$r(x)=\underset{k=0}{\overset{\mathrm{}}{}}\frac{\mathrm{\Gamma }(k\gamma +1/4i\gamma \stackrel{~}{\theta }/2\pi )\mathrm{\Gamma }((k+1)\gamma +3/4i\gamma \stackrel{~}{\theta }/2\pi )}{\mathrm{\Gamma }((k+1/2)\gamma +1/4i\gamma \stackrel{~}{\theta }/2\pi )\mathrm{\Gamma }((k+1/2)\gamma +3/4i\gamma \stackrel{~}{\theta }/2\pi )},$$
(3.44)
and $`\stackrel{~}{\theta }=\theta \eta `$. Therefore,
$$A(\theta )=\frac{1}{\sqrt{\nu }}\frac{e^{i\pi (1+\gamma )/4}}{(1+ipx)}\frac{r(x)}{\overline{r}(x)}.$$
(3.45)
If desired, this may be written in a more compact form obtained by making repeated use of Legendre’s formula
$$\mathrm{\Gamma }(z)=\frac{2^{z1/2}}{\sqrt{2\pi }}\mathrm{\Gamma }\left(\frac{z}{2}\right)\mathrm{\Gamma }\left(\frac{z}{2}+\frac{1}{2}\right).$$
(3.46)
Thus, for example,
$$\frac{r(x)}{\overline{r}(x)}=\frac{\mathrm{\Gamma }(1/4i\gamma \stackrel{~}{\theta }/2\pi )}{\mathrm{\Gamma }(1/4+i\gamma \stackrel{~}{\theta }/2\pi )}\underset{k=1}{\overset{\mathrm{}}{}}\frac{\mathrm{\Gamma }(1/2+2k\gamma i\gamma \stackrel{~}{\theta }/\pi )\mathrm{\Gamma }(1/2+(2k1)\gamma +i\gamma \stackrel{~}{\theta }/\pi )}{\mathrm{\Gamma }(1/2+2k\gamma +i\gamma \stackrel{~}{\theta }/\pi )\mathrm{\Gamma }(1/2+(2k1)\gamma i\gamma \stackrel{~}{\theta }/\pi )}.$$
(3.47)
It is clear from what has been said so far that the above arguments could be repeated for a transmission matrix with odd defect labels. However, to do so might be misleading and in particular the unitarity relations are inappropriate since it is expected, as already mentioned, that the odd labelled defect is unstable.
Since the parameter $`p=|\mu |=e^{\gamma \eta }`$ will turn out to be important it is worth resummarising the solution for the components of $`{}_{}{}^{\mathrm{e}}T`$:
$`A_\alpha ^\beta =Q^\alpha A(\theta )\delta _\alpha ^\beta ,`$ $`D_\alpha ^\beta =\nu Q^\alpha A(\theta )\delta _\alpha ^\beta ,`$ (3.48)
$`B_\alpha ^\beta =\stackrel{~}{\lambda }\sigma ^\alpha px_\theta A(\theta )\delta _\alpha ^{\beta 2},`$ $`C_\alpha ^\beta =\stackrel{~}{\mu }\sigma ^\alpha px_\theta A(\theta )\delta _\alpha ^{\beta +2},`$
where $`\mu =p\stackrel{~}{\mu }`$ and $`\stackrel{~}{\lambda }=\nu \sigma ^2/q\stackrel{~}{\mu }`$. The phases $`\nu ,\stackrel{~}{\mu },\sigma `$ are all independent of $`\theta `$ and, if they are also independent of $`\eta `$, the transmission matrix turns out to be a function only of the difference $`\theta \eta `$.
One further remark before proceeding. It is possible to make a change of basis, redefining the transmission matrix via a diagonal unitary transformation without upsetting the assumptions made at the beginning of the section. This must be a transformation on the defect labels alone, and it may be written in the form
$$V_\alpha ^\beta =v_\alpha \delta _\alpha ^\beta ,|v_\alpha |=1,$$
(3.49)
for which $`A`$ and $`D`$ are unchanged, but $`B`$ and $`C`$ alter according to,
$$b_\alpha v_\alpha b_\alpha \overline{v}_{\alpha +2},c_\alpha v_\alpha c_\alpha \overline{v}_{\alpha 2}.$$
(3.50)
Such a transformation can be designed to remove the phases $`\stackrel{~}{\lambda }\sigma ^\alpha `$ from $`b_\alpha `$ by selecting
$$v_\alpha =\sigma ^{\alpha (\alpha 2)/4}\stackrel{~}{\lambda }^{\alpha /2},$$
and then
$$c_\alpha \stackrel{~}{\mu }\stackrel{~}{\lambda }\sigma ^2=\nu /q.$$
A further transformation of the same type can then be made to ensure the coefficients in the off-diagonal terms are the same. In other words, the dependence on $`\sigma `$ and the two phases $`\stackrel{~}{\mu },\stackrel{~}{\lambda }`$ can be removed entirely to leave the symmetrical expression
$`\left(\begin{array}{cc}A_\alpha ^\beta & B_\alpha ^\beta \\ C_\alpha ^\beta & D_\alpha ^\beta \end{array}\right)=f(q,x)\left(\begin{array}{cc}\nu ^{1/2}Q^\alpha \delta _\alpha ^\beta & q^{1/2}e^{\gamma (\theta \eta )}\delta _\alpha ^{\beta 2}\\ q^{1/2}e^{\gamma (\theta \eta )}\delta _\alpha ^{\beta +2}& \nu ^{1/2}Q^\alpha \delta _\alpha ^\beta \end{array}\right)`$ (3.55)
The result (3.55) has been arrived at via a collection of arguments based on the triangle relation (3.9) and general principles. Before proceeding to investigate some of its consequences, it is worth noting some general features and comparing the result with what one might have expected on the basis of the Lagrangian description. If $`\eta >0`$ and $`\theta >\eta `$, the character-changing elements dominate while for $`\theta <\eta `$ they are suppressed. On the other hand, for $`\eta <0`$, the character-changing processes are always dominant. All these facts are compatible with the classical transmission properties described in section (2). Also, it is striking that the matrix elements representing a soliton converting to an anti-soliton or vice-versa are the same in this basis, while elements representing the scattering of solitons or anti-solitons with no change of character are different. This aspect follows from an argument based on the Lagrangian and makes use of the particular type of defect described there.
Suppose the transmission matrix is expressed formally as a functional integral over the fields $`u`$ and $`v`$, weighted by the classical action including the defect contribution (2.2); and further suppose that a defect on its own is labelled by the vacuum configurations of the fields to either side of it. Thus the label $`(a,b)`$ is ascribed to the defect when the fields have the constant values $`u=2a\pi /\beta ,v=2b\pi /\beta `$. Provided the labels range over even integers the defect potential and the bulk terms have the values they would have for $`(a,b)=(0,0)`$, corresponding to a stable defect ($`\sigma >0)`$. The topological charge of the defect is $`ba`$. The term in the Lagrangian linear in the time derivatives has a consequence that may be explored as follows. Field configurations $`u,v`$ evolving in the presence of an initial defect with labels $`(a,b)`$ may be compared with configurations evolving with an initial defect labelled $`(0,0)`$ by translating the fields via $`uu2a\pi /\beta ,vv2b\pi /\beta `$. The bulk and defect potential parts do not change under this but, due to the terms linear in time derivatives, the action changes by a term
$$\frac{\pi }{\beta }_{\mathrm{}}^{\mathrm{}}𝑑t\left(av_tbu_t\right)_{x=0}=\frac{\pi }{\beta }(a\delta vb\delta u),$$
(3.56)
where $`\delta u,\delta v`$ are the nett changes over time of the field configurations evaluated at the location of the defect. Consequently, the functional integral representing the transmission factor relative to an initial defect $`(a,b)`$ will differ from the transmission relative to $`(0,0)`$ by a constant factor given by
$$T(a,b)=\mathrm{exp}\left(\frac{i\pi }{\beta }(a\delta vb\delta u)\right)T(0,0).$$
(3.57)
For example, a soliton travelling in a positive sense along the $`x`$axis will produce the shifts $`(a,b)(a1,b1),\delta u=2\pi /\beta =\delta v`$, if its character does not change, thus acquiring a factor
$$e^{2i\pi ^2(ba)/\beta ^2}Q^{\frac{(ba)}{2}}.$$
(3.58)
On the other hand, it requires the shifts $`(a,b)(a1,b+1),\delta u=2\pi /\beta =\delta v`$ and acquires a factor
$$e^{2i\pi ^2(b+a)/\beta ^2}Q^{\frac{(b+a)}{2}},$$
(3.59)
if its character changes. The corresponding factors for an anti-soliton are, respectively,
$$Q^{\frac{(ba)}{2}}\text{and}Q^{\frac{(b+a)}{2}}.$$
(3.60)
Previously, the defect was labelled by $`\alpha =ba`$, its topological charge, but the above remarks suggest there might also be a dependence on the quantity $`p=a+b`$. Thus, in the notation introduced at the beginning of this section, in (3.11), for the components of $`{}_{}{}^{\mathrm{e}}T`$, it is tempting to write
$`A_{\alpha p}^{\beta q}=Q^{\alpha /2}\widehat{a}_0\delta _\alpha ^\beta \delta _p^{q+2},`$ $`B_{\alpha p}^{\beta q}=Q^{p/2}\widehat{b}_0\delta _\alpha ^{\beta 2}\delta _p^q,`$
$`C_{\alpha p}^{\beta q}=Q^{p/2}\widehat{c}_0\delta _\alpha ^{\beta +2}\delta _p^q,`$ $`D_{\alpha p}^{\beta q}=Q^{\alpha /2}\widehat{d}_0\delta _\alpha ^\beta \delta _p^{q2}.`$ (3.61)
The expressions (3) appear to have a different dependence on $`Q`$ to that reported in (3.55) and include coupling-dependent terms sensitive to $`p=a+b`$. However, the dependence on $`Q`$ may be adjusted using the diagonal unitary transformation
$$V_{\alpha p}^{\beta q}=Q^{p\alpha /4}\delta _\alpha ^\beta \delta _p^q,$$
(3.62)
to obtain instead
$`A_{\alpha p}^{\beta q}=Q^\alpha \widehat{a}_0\delta _\alpha ^\beta \delta _p^{q+2},`$ $`B_{\alpha p}^{\beta q}=\widehat{b}_0\delta _\alpha ^{\beta 2}\delta _p^q,`$
$`C_{\alpha p}^{\beta q}=\widehat{c}_0\delta _\alpha ^{\beta +2}\delta _p^q,`$ $`D_{\alpha p}^{\beta q}=Q^\alpha \widehat{d}_0\delta _\alpha ^\beta \delta _p^{q2}.`$ (3.63)
Using this as a starting point and applying the same arguments as before leads to a similar expression to (3.55), the only changes being the extra Kronecker $`\delta `$’s needed to record the changes in $`a+b`$ as a soliton or anti-soliton passes the defect. In fact, a unitary basis can be found which removes from this representation even the dependence on $`\nu `$, the undetermined parameter in (3.55). On the other hand, although the transmission matrix has not been defined explicitly by a functional integral, the arguments given above confirm both the importance of the term linear in time derivatives and the precise dependence on the coupling that has been coded into $`Q`$. As a consequence, if it ever becomes possible to give an independent derivation of the transmission matrix then the triangle relations (3.9) could be used to provide an alternative derivation of the soliton S-matrix.
## 4 Unstable bound states as poles in $`{}_{}{}^{\mathrm{e}}T`$
Returning to an earlier line of thought, consider the poles of $`{}_{}{}^{\mathrm{e}}T(\theta )`$. At first sight it seems from (3.43) that each component of $`{}_{}{}^{\mathrm{e}}T(\theta )`$ has a pole at $`px=i`$, or in terms of rapidity at $`\theta =\eta +i\pi /2\gamma `$ (or, equivalently, $`i\gamma \stackrel{~}{\theta }/2\pi =1/4`$). However, this is an illusion since $`\overline{r}(x)`$ also has a pole at the same location which cancels it out. On the other hand, it is evident from the expressions (3.47) and (3.40) that there is a pole at $`\theta =\eta i\pi /2\gamma `$ (that is, $`i\gamma \stackrel{~}{\theta }/\pi =1/2`$), and this pole does not have a compensating zero. If this pole is taken to be significant then it corresponds to a defect ‘bound’ state with an energy jump given by
$$\mathrm{\Delta }E=E_bE_0=m_s\mathrm{cosh}\left(\eta \frac{i\pi }{2\gamma }\right)=m_s\mathrm{cosh}\eta \mathrm{cos}\left(\frac{\pi }{2\gamma }\right)im_s\mathrm{sinh}\eta \mathrm{sin}\left(\frac{\pi }{2\gamma }\right),$$
(4.1)
where $`m_s`$ is the mass of the soliton or anti-soliton. Equation (4.1) is interesting because in view of (3.8), the limit $`\beta 0`$ implies $`1/\gamma 0`$ and the real part of $`\mathrm{\Delta }E`$ becomes the classical energy of a soliton with rapidity $`\eta `$, and the imaginary part of $`\mathrm{\Delta }E`$ vanishes. This suggests this pole corresponds to the absorption and emission of a soliton (or anti-soliton) and the imaginary part of the pole location governs the width of the ‘bound’ state (in other words its decay, corresponding to emission). The sign of the imaginary part is reasonable since one would expect the wave-function of the bound state to evolve according to
$$\psi _b(t)=e^{iE_bt}\psi _b(0).$$
(4.2)
Thus, $`\psi _b(t)`$ is certainly decaying with time provided the imaginary part of $`\mathrm{\Delta }E`$ is negative. In the classical limit the width goes to zero so the only process is absorption. From this point of view, it is fortunate the other pole turned out to be a phantom.
Another interesting point concerns the picture if momentum is incorporated. Since a particle generally passes through a defect, and in the classical picture can exchange both energy and momentum with it, it becomes clear from the expressions for the energy and momentum that the energy-momentum of the excited quantum defect is a complex rotation of the energy momentum of a bulk soliton with rapidity $`u`$. Thus,
$$\left(\begin{array}{c}\mathrm{\Delta }E\\ \mathrm{\Delta }P\end{array}\right)=\left(\begin{array}{cc}\mathrm{cos}\chi & i\mathrm{sin}\chi \\ i\mathrm{sin}\chi & \mathrm{cos}\chi \end{array}\right)\left(\begin{array}{c}m_s\mathrm{cosh}\eta \\ m_s\mathrm{sinh}\eta \end{array}\right),\chi =\frac{\pi }{2\gamma }.$$
(4.3)
Clearly,
$$(\mathrm{\Delta }E)^2(\mathrm{\Delta }P)^2=m_s^2,$$
(4.4)
and hence, despite having a complex energy and momentum suited to representing an unstable state, the excited state remains on the soliton mass shell. This seems to fit exceedingly well with the classical picture of a defect (of odd topological charge and therefore additional energy) being a ‘hidden’ soliton.
Another interesting limit corresponds to $`\beta ^2=4\pi `$ ($`\gamma =1`$). In that limit, the solitons become free , and the real parts of $`\mathrm{\Delta }E`$ and $`\mathrm{\Delta }P`$ vanish. Nevertheless, the bound state remains on the soliton mass shell. In that limit, the rôles of energy and momentum have completely interchanged and, in that sense, one might wonder about an analogy with a black hole described by the Schwarzschild metric for which ‘timelike’ and ‘spacelike’ interchange roles at the horizon. In this limit, $`Q^2=1`$ and the transmission factor (3.55) simplifies considerably to
$$^\mathrm{e}T_{a\alpha }^{b\beta }(\theta )=\frac{1}{(1ie^{\stackrel{~}{\theta }})}\left(\begin{array}{cc}\nu ^{1/2}\delta _\alpha ^\beta & ie^{\stackrel{~}{\theta }}\delta _\alpha ^{\beta 2}\\ ie^{\stackrel{~}{\theta }}\delta _\alpha ^{\beta +2}& \nu ^{1/2}\delta _\alpha ^\beta \end{array}\right).$$
(4.5)
Near the bound state pole, the transmission matrix should have the form
$$^\mathrm{e}T_{a\alpha }^{b\beta }(\theta )\frac{it_{a\alpha }^{b\beta }}{\left(\theta \eta +\frac{i\pi }{2\gamma }\right)}\frac{ic_{a\alpha }^\delta \stackrel{~}{c}_\delta ^{b\beta }}{\left(\theta \eta +\frac{i\pi }{2\gamma }\right)},$$
(4.6)
where $`c_{a\alpha }^\delta `$ and $`\stackrel{~}{c}_\delta ^{b\beta }`$ are a pair of coupling matrices factorizing $`t_{a\alpha }^{b\beta }`$. Note that while $`\alpha `$ and $`\beta `$ are even, the internal sum implied by the repeated $`\delta `$ is over the odd integers. However, there is clearly considerable freedom in choosing these couplings. An explicit calculation of the pole residue reveals - in the notation of (3.55),
$$^\mathrm{e}T_{a\alpha }^{b\beta }(\theta )\frac{1}{\gamma }\frac{e^{i\pi (1+\gamma )/4}}{(\theta \eta +\frac{i\pi }{2\gamma })}\left(\begin{array}{cc}i\nu ^{1/2}Q^\alpha \delta _\alpha ^\beta & q^{1/2}\delta _\alpha ^{\beta 2}\\ q^{1/2}\delta _\alpha ^{\beta +2}& i\nu ^{1/2}Q^\alpha \delta _\alpha ^\beta \end{array}\right).$$
(4.7)
Some care needs to be taken when manipulating the infinite products of gamma functions but the result (4.7) can be checked with the special case $`\gamma =1`$.
The intermediate state, of complex energy $`E_b`$, if this is a correct interpretation, should differ from the initial state of energy $`E_0`$ by a single unit of topological charge. For this reason, there will be a bootstrap condition linking the even and odd transmission matrices. In detail, it ought to read
$$c_{b\alpha }^\gamma {}_{}{}^{\mathrm{o}}T_{a\gamma }^{c\delta }=S_{ab}^{pq}\left(\theta \eta +\frac{i\pi }{2\gamma }\right)^\mathrm{e}T_{q\alpha }^{c\beta }c_{p\beta }^\delta ,$$
(4.8)
where, for example, $`\alpha ,\beta `$ are both even, and $`\gamma ,\delta `$ are both odd. One can argue to this assignment by considering the classical energy of the defect poised to emit a soliton: the odd-charged defect has higher energy and hence it is to be expected that the above pole really does refer to a pole in the transmission matrix $`{}_{}{}^{\mathrm{e}}T`$ for a particle passing through an even-charged defect. Then it is clear $`{}_{}{}^{\mathrm{o}}T`$ refers to a particle scattering with an odd-charge defect. It is worth pointing out that the bootstrap relation and the Yang-Baxter equation satisfied by the $`S`$-matrix, along with (3.9), will guarantee that $`{}_{}{}^{\mathrm{o}}T`$ satisfies a corresponding set of equations to (3.9), and therefore the general solution will be of a similar form to (3). The corresponding quantities will be denoted by $`\widehat{A}(\theta ),\widehat{\nu },`$ etc. However, although the equation corresponding to (3.35) is expected to be satisfied, it is not expected there will be an analogue of the unitarity condition (3.1).
In order to facilitate the checking of (4.8) it will be necessary to calculate the couplings $`c_{a\alpha }^\beta `$ and to evaluate the sine-Gordon S-matrix at a special point. The simplest example of (4.8) concerns a soliton scattering with a soliton-defect bound state since then one must have,
$$c_{+\alpha }^\gamma {}_{}{}^{\mathrm{o}}T_{+\gamma }^{+\delta }(\theta )=S_{++}^{++}\left(\theta \eta +\frac{i\pi }{2\gamma }\right)^\mathrm{e}T_{+\alpha }^{+\beta }c_{+\beta }^\delta ,$$
(4.9)
which in terms of the specific entries of the transmission matrices translates into
$$c_{+\alpha }^\gamma \widehat{a}_\gamma =S_{++}^{++}\left(\theta \eta +\frac{i\pi }{2\gamma }\right)a_\alpha c_{+\alpha }^\gamma .$$
(4.10)
However, on the left hand side it is clear $`\gamma =\alpha +1`$ and so the couplings cancel (assuming they are invertible matrices) and
$$Q^{\alpha +1}\widehat{A}(\theta )=S_{++}^{++}\left(\theta \eta +\frac{i\pi }{2\gamma }\right)Q^\alpha A(\theta ).$$
(4.11)
Since, the right hand side is in principle known, the left hand side can be calculated. If this idea is to be consistent there must be some very interesting identities involving the S-matrix elements evaluated at a rapidity shifted by the position of the bound state pole. For example, it needs to be checked that this relation is consistent with what has been assumed already concerning the transmission matrix. The first useful identity satisfied by the S-matrix is
$$S_{++}^{++}(\theta \eta +i\psi )S_{++}^{++}(\theta \eta +i\psi +i\pi )=q\frac{1e^{2\gamma (i\psi \eta )}x^2}{1e^{2\gamma (i\psi \eta )}q^2x^2},$$
(4.12)
where $`\eta ,\psi `$ are real. On the other hand, (4.11) implies
$$Q^2\widehat{A}(\theta )\widehat{A}(\theta +i\pi )=S_{++}^{++}(\theta \eta +i\psi )S_{++}^{++}(\theta \eta +i\psi +i\pi )A(\theta )A(\theta +i\pi ),$$
(4.13)
and hence, using (3.35), one finds
$$\frac{1}{\widehat{\nu }(1+\widehat{p}^2x^2)}=\left(\frac{1e^{2\gamma (i\psi \eta )}x^2}{1e^{2\gamma (i\psi \eta )}q^2x^2}\right)\frac{1}{\nu (1+p^2x^2)}.$$
(4.14)
Since $`x=e^{\gamma \theta }`$ is free, comparing powers of $`x`$ leads to the conclusions:
$$\widehat{\nu }=\nu ,\widehat{p}^2=q^2p^2,\psi =\frac{\pi }{2\gamma }\pm \frac{k\pi }{\gamma },k\mathrm{integer}.$$
(4.15)
An explicit computation of $`\widehat{A}(\theta )`$ using (4.11) yields,
$$\widehat{A}(\theta )=\frac{1}{\sqrt{\nu }}\frac{e^{i\pi (1+\gamma )/4}}{(1+ipx)}\frac{\mathrm{cos}\left(\pi /4\gamma i(\theta \eta )/2\right)}{\mathrm{sin}\left(\pi /4\gamma i(\theta \eta )/2\right)}\frac{s(x)}{\overline{s}(x)},$$
(4.16)
where $`s(x)`$ is quite similar to $`r(x)`$ defined in (3.44)
$$s(x)=\underset{k=0}{\overset{\mathrm{}}{}}\frac{\mathrm{\Gamma }(k\gamma +3/4+i\gamma \stackrel{~}{\theta }/2\pi )\mathrm{\Gamma }((k+1)\gamma +1/4+i\gamma \stackrel{~}{\theta }/2\pi )}{\mathrm{\Gamma }((k+1/2)\gamma +1/4+i\gamma \stackrel{~}{\theta }/2\pi )\mathrm{\Gamma }((k+1/2)\gamma +3/4+i\gamma \stackrel{~}{\theta }/2\pi )}.$$
(4.17)
Using (4.16) it is not difficult to check directly that the condition (3.35) holds but with $`p`$ replaced by $`\widehat{p}`$, namely
$$\widehat{A}(\theta )\widehat{A}(\theta +i\pi )=\frac{1}{\nu (1+\widehat{p}^2x^2)}.$$
Note, no further information is obtained from considering an anti-soliton scattering with an anti-soliton-defect bound state since that would require
$$c_\alpha ^\gamma \widehat{d}_\gamma =S_{}^{}\left(\theta \eta +\frac{i\pi }{2\gamma }\right)d_\alpha c_\alpha ^\gamma .$$
(4.18)
On the left hand side $`\gamma =\alpha 1`$ and therefore, recalling the properties of the S-matrix and that $`d_\alpha =\nu Q^\alpha A(\theta )`$, one finds
$$Q\widehat{\nu }\widehat{A}(\theta )=S_{++}^{++}\left(\theta \eta +\frac{i\pi }{2\gamma }\right)\nu A(\theta ).$$
(4.19)
This is exactly the previous relation because $`\widehat{\nu }=\nu `$. Clearly, from the relations (4.15), when there is a complex bound state pole not both of $`p`$ and $`\widehat{p}`$ can be real. In turn this confirms that the unitarity assumption must be faulty for $`{}_{}{}^{\mathrm{o}}T(\theta )`$, as expected. On the other hand, the relationship corresponding to (3.35) survives. In fact, the incompatibility of (4.11) with unitarity can be checked directly by multiplying each side of (4.11) by its complex conjugate and comparing with (3.37). The fact the S-matrix is evaluated at $`\theta \eta +i\pi /2\gamma `$ provides an immediate contradiction since
$$\overline{S}_{++}^{++}(\theta \eta +i\pi /2\gamma )S_{++}^{++}(\theta \eta +i\pi /2\gamma )1.$$
To obtain further information from (4.8) would require more complete knowledge of the coupling matrices.
## 5 Transmission factors for the breathers
In this section the precise form of the transmission factor for the lightest breather will be calculated and it turns out to be surprising. Knowledge of this particular transmission factor should provide a way to begin checking perturbatively the properties of the defect: since the lightest breather corresponds to the quantum particle described by the fundamental bulk scalar field, a development of standard methods can be used to examine low orders of perturbation theory.
The breather bound state poles of the soliton-soliton scattering matrix (3.2) occur at
$$\mathrm{\Theta }=i(\pi \frac{n\pi }{\gamma }),n=1,2,3,\mathrm{}.,n_{\mathrm{max}},$$
(5.1)
where $`n_{\mathrm{max}}`$ depends upon the coupling $`\beta `$ and is defined as the largest integer for which $`n/\gamma <1`$. The corresponding breather masses are given by
$$m_n=2m_s\mathrm{sin}\left(\frac{n\pi }{2\gamma }\right).$$
(5.2)
The existence of the breather poles allows a calculation of breather transmission factors $`{}_{}{}^{n}T(\theta )`$ via a bootstrap relation as follows
$$c_{a\overline{a}}^n{}_{}{}^{n}T(\theta )\delta _\alpha ^\beta =\underset{b}{\overset{\mathrm{e}}{}}T_{\overline{a}\alpha }^{\overline{b}\gamma }(\theta _{\overline{a}})^\mathrm{e}T_{a\gamma }^{b\beta }(\theta _a)c_{b\overline{b}}^n,$$
(5.3)
where
$$\theta _a=\theta +i\left(\frac{\pi }{2}\frac{n\pi }{2\gamma }\right),\theta _{\overline{a}}=\theta i\left(\frac{\pi }{2}\frac{n\pi }{2\gamma }\right),$$
(5.4)
and the label $`b`$ is summed over the two possibilities $`b=\pm `$. It is necessary to be slightly careful with the couplings in (5.3) because
$$c_+^n=()^nc_+^n,$$
especially since the main interest here lies with the first breather $`(n=1)`$. Taking that into account and using the explicit expressions for the soliton transmission matrix elements given in (3) the breather transmission matrix for the lightest breather is given by
$$^1T(\theta )\delta _\alpha ^\beta =(D(\theta _{\overline{a}})A(\theta _a)C(\theta _{\overline{a}})B(\theta _a))_\alpha ^\beta =(1\frac{p^2x^2}{q})\nu A(\theta _{\overline{a}})A(\theta _a)\delta _\alpha ^\beta .$$
(5.5)
It is straightforward to use the expression for $`A(\theta )`$ provided by (3.43) to calculate $`{}_{}{}^{1}T(\theta )`$, the only additional formula needed being the infinite product identity for the sine function
$$\frac{\mathrm{sin}\pi z}{\pi z}=\underset{k=1}{\overset{\mathrm{}}{}}\left(1\frac{z^2}{k^2}\right).$$
(5.6)
The result is surprising since all dependence on the bulk coupling cancels out to leave:
$$^1T(\theta )=i\frac{\mathrm{sin}\left(\frac{\pi }{4}+\frac{i(\theta \eta )}{2}\right)}{\mathrm{sin}\left(\frac{\pi }{4}\frac{i(\theta \eta )}{2}\right)}i\frac{\mathrm{sinh}\left(\frac{\theta \eta }{2}\frac{i\pi }{4}\right)}{\mathrm{sinh}\left(\frac{\theta \eta }{2}+\frac{i\pi }{4}\right)}.$$
(5.7)
Moreover, this expression is identical to the transmission factor in the classical (linear) limit given in section (2) by (2.9). Given the delicate cancellations this seems a noteworthy result.
Given that the sinh-Gordon model is the imaginary coupling version of the sine-Gordon model, the result (5.7) also suggests that the scalar particle of the sinh-Gordon model has a transmission factor independent of the bulk coupling.
Although the result (5.7) is surprising, it might be altered if the ‘minimal’ version of $`{}_{}{}^{\mathrm{e}}T(\theta )`$ were to be modified by the addition of CDD factors. The purpose of a perturbative analysis would be to establish the validity, or otherwise, of (5.7), and hence to determine whether or not CDD factors might be necessary. Resolving this issue will be postponed to a future report.
The transmission factors for the other breathers may be obtained iteratively using the bootstrap. For example, denoting the scattering matrix of the lightest breather by $`S_{11}(\mathrm{\Theta })`$,
$$S_{11}(\mathrm{\Theta })=\frac{\mathrm{sinh}\left(\frac{\mathrm{\Theta }}{2}+\frac{i\pi }{2\gamma }\right)\mathrm{sinh}\left(\frac{\mathrm{\Theta }}{2}\frac{i\pi }{2}\left(1+\frac{1}{\gamma }\right)\right)}{\mathrm{sinh}\left(\frac{\mathrm{\Theta }}{2}\frac{i\pi }{2\gamma }\right)\mathrm{sinh}\left(\frac{\mathrm{\Theta }}{2}+\frac{i\pi }{2}\left(1+\frac{1}{\gamma }\right)\right)}=(2/\gamma )_\mathrm{\Theta }(22/\gamma )_\mathrm{\Theta },$$
(5.8)
where the last expression makes use of the bracket notation
$$(z)_\theta =\frac{\mathrm{sinh}\left(\frac{\theta }{2}+\frac{i\pi z}{4}\right)}{\mathrm{sinh}\left(\frac{\theta }{2}\frac{i\pi z}{4}\right)},$$
(5.9)
there is a pole corresponding to the next breather at $`\mathrm{\Theta }=i\pi /\gamma `$. Using this the transmission factor of the second breather will be given by
$$^2T(\theta )=^1T(\stackrel{~}{\theta }+\frac{i\pi }{2\gamma })^1T(\stackrel{~}{\theta }\frac{i\pi }{2\gamma })=\frac{1}{(1+1/\gamma )_{\stackrel{~}{\theta }}(11/\gamma )_{\stackrel{~}{\theta }}},$$
(5.10)
where $`\stackrel{~}{\theta }=\theta \eta `$. The result is obviously dependent on the bulk coupling. Unlike the expression (5.7) the transmission factor for the second breather contains two complex poles possibly indicating the existence of an excited state of the defect with no change in its topological charge. However, using the same argument as before, any defect-breather bound state would remain on the mass shell of the breather and would need to be interpreted as the quantised version of a phenomenon in the classical model. However, mysteriously, there do not appear to be any periodic solutions associated specifically with the defect. On the other hand, from an energy relation similar to (4.1) it is clear that for the two possible poles, $`\stackrel{~}{\theta }=i\pi (1\pm 1/\gamma )/2`$,
$$\mathrm{\Delta }E_\pm =m_2\mathrm{cosh}\eta \mathrm{sin}\left(\frac{\pi }{2\gamma }\right)im_2\mathrm{sinh}\eta \mathrm{cos}\left(\frac{\pi }{2\gamma }\right),$$
(5.11)
while the pole coefficients are $`\pm 2i\mathrm{cot}(\pi /2\gamma )`$, respectively. Thus, the pole with a positive residue has $`\mathrm{Re}\mathrm{\Delta }E<0`$, and the pole with $`\mathrm{Re}\mathrm{\Delta }E>0`$ has a negative residue. As $`\beta 0`$ (i.e. $`1/\gamma 0`$), the pair of poles coalesce to a double pole. For $`\gamma >1`$, neither of these poles is on the ‘physical strip’ ($`0<\mathrm{Im}\stackrel{~}{\theta }<\pi `$). Nevertheless, it is an interesting question to decide the nature of the states indicated by these poles, or what their origin might be if they do not correspond to unstable states.
The general case needs to be split into ‘odd’ and ‘even’ breathers, according to whether $`n`$ in (5.1) is odd or even. In summary, repeatedly using (5.3) one finds,
$$^{2\mathrm{s}+1}T(\theta )=()^{s+1}\frac{i}{(1)_{\stackrel{~}{\theta }}}\underset{l=0}{\overset{s1}{}}\frac{1}{(1+2(sl)/\gamma )_{\stackrel{~}{\theta }}(12(sl)/\gamma )_{\stackrel{~}{\theta }}},$$
(5.12)
for $`n`$ odd, and
$$^{2\mathrm{s}}T(\theta )=()^s\underset{l=0}{\overset{s1}{}}\frac{1}{(1+(2s2l1)/\gamma )_{\stackrel{~}{\theta }}(1(2s2l1)/\gamma )_{\stackrel{~}{\theta }}},$$
(5.13)
for $`n`$ even. Typically, in common with the transmission factor for the second breather, these have complex poles suggesting the existence of additional excited states.
## 6 Scattering defects
As remarked earlier in section (3) there is no reason in principle why defects themselves will not scatter and the purpose of this section is to make some comments about this. Throughout the section it will be assumed the defect parameter is positive, so that the even type of defect is stable. Then, it could be expected that there would be asymptotic states composed of any number of defects with different speeds. It might be expected the speed of a defect will modify the defect parameter but it is necessary to redo some classical parts of the analysis to determine if this is really correct and if so, how the transmission matrix already determined will be modified. It was noted in how the requirement of having a conserved total momentum, including a defect contribution, was equivalent to the requirements of integrability. In fact, the existence of the conserved momentum implied the form of the defect boundary conditions. This observation will be used as a shortcut in the present discussion.
Starting from first principles, the action with a moving defect is taken to be
$$A=𝑑t\left\{_{\mathrm{}}^z𝑑x(u)+𝐁m_D\sqrt{1\dot{z}^2}+_z^{\mathrm{}}𝑑x(v)\right\},$$
(6.1)
where $`z`$ depends on time and $`𝐁`$ depends on the fields evaluated at $`x=z(t)`$, and possibly on $`z`$ or $`\dot{z}`$ in a manner to be determined. A velocity dependent action for the defect is included assuming its mass to be $`m_D`$. On the defect, the total time derivative of either field (or its variation) will be given in terms of a combination such as
$$\dot{u}=\frac{du}{dt}=\frac{u}{t}+\dot{z}\frac{u}{x}u_t+\dot{z}u_x,$$
(6.2)
or a similar expression for $`\dot{v}`$.
Then, requiring $`A`$ to be stationary with respect to variations of $`u`$ gives:
$`{\displaystyle \frac{}{u_x}}+{\displaystyle \frac{𝐁}{u}}{\displaystyle \frac{d}{dt}}{\displaystyle \frac{𝐁}{u_t}}\dot{z}{\displaystyle \frac{}{u_t}}`$ $`=`$ $`0,x=z`$ (6.3)
$`{\displaystyle \frac{𝐁}{u_x}}\dot{z}{\displaystyle \frac{𝐁}{u_t}}`$ $`=`$ $`0,x=z`$ (6.4)
$`{\displaystyle \frac{}{u}}_\mu {\displaystyle \frac{}{u_\mu }}`$ $`=`$ $`0,x<z.`$ (6.5)
Equation (6.4) simply implies $`𝐁`$ depends on the combination $`du/dt`$ while equation (6.3) gives the boundary condition for $`u/x`$:
$$u_x=\frac{𝐁}{u}\frac{d}{dt}\frac{𝐁}{u_t}\dot{z}u_t,x=z;$$
(6.6)
this suggests it is natural to choose as a defect contribution
$$𝐁=\frac{1}{2}\left(u\frac{dv}{dt}v\frac{du}{dt}\right).$$
(6.7)
In the latter, it may then be assumed $``$ will have no dependence on the derivatives of the fields (but will still depend on $`z,\dot{z}`$).
Then, (and similarly for $`v`$),
$$u_x=\frac{1}{1\dot{z}^2}\left(\frac{}{u}+\frac{dv}{dt}\dot{z}\frac{du}{dt}\right),v_x=\frac{1}{1\dot{z}^2}\left(\frac{}{v}+\frac{du}{dt}\dot{z}\frac{dv}{dt}\right),$$
(6.8)
both of which reduce to the old situation when $`z`$ is constant. Alternatively, one can combine these differently to obtain:
$$u_x=\frac{1}{1\dot{z}^2}\left(\frac{}{u}+\dot{z}\frac{}{v}\right)+v_t,v_x=\frac{1}{1\dot{z}^2}\left(\frac{}{v}\dot{z}\frac{}{u}\right)+u_t.$$
(6.9)
For some purposes the latter is simpler. There are similar expressions for the partial time derivatives at the defect:
$$u_t=\frac{1}{1\dot{z}^2}\left(\dot{z}\frac{}{u}+\frac{du}{dt}\dot{z}\frac{dv}{dt}\right),v_t=\frac{1}{1\dot{z}^2}\left(\dot{z}\frac{}{v}+\frac{dv}{dt}\dot{z}\frac{du}{dt}\right);$$
(6.10)
these are identities when $`z`$ is constant.
Next, consider the contribution of the two bulk fields $`u`$ and $`v`$ to the total momentum,
$$P=_{\mathrm{}}^z𝑑xu_tu_x+_z^{\mathrm{}}𝑑xv_tv_x.$$
(6.11)
This is not expected to be conserved. However, the time derivative of $`P`$ is
$`\dot{P}`$ $`=`$ $`\left[\dot{z}\left(u_tu_xv_tv_x\right)+{\displaystyle \frac{1}{2}}\left(u_t^2+u_x^2v_t^2v_x^2\right)\left(V(u)W(v)\right)\right]_{x=z},`$ (6.12)
and the various partial derivatives evaluated at the defect can all be related to the total time derivatives using the defect conditions for $`u_x,v_x`$, and the definition of the total time derivative. One finds that all terms quadratic in total time derivatives cancel out; that terms linear in total derivatives reduce to
$$\frac{1}{1\dot{z}^2}\left[\dot{u}\left(\dot{z}\frac{}{u}\frac{}{v}\right)+\dot{v}\left(\dot{z}\frac{}{v}\frac{}{u}\right)\right];$$
(6.13)
and that terms without any derivatives at all simplify to
$$\frac{1}{2(1\dot{z}^2)}\left[\left(\frac{}{u}\right)^2\left(\frac{}{v}\right)^2\right]V(u)+W(v).$$
(6.14)
Consider first the last of these. The two potentials do not depend on $`\dot{z}`$ and, therefore, a suitable proposal for $``$ will be
$$=\sqrt{1\dot{z}^2}𝒞,$$
(6.15)
where $`𝒞`$ is chosen to ensure (6.14) vanishes identically when $`\dot{z}=0`$. Assuming this to be the case, (6.13) is quite close to a total time derivative. Setting
$$\frac{𝒞}{u}=\frac{𝒰}{v},\frac{𝒞}{v}=\frac{𝒰}{u},$$
(6.16)
(6.13) becomes
$$\frac{1}{\sqrt{1\dot{z}^2}}\left[\dot{u}(\dot{z}𝒞_u𝒰_u)+\dot{v}(\dot{z}𝒞_v𝒰_v)\right].$$
(6.17)
Finally, if $`𝒞`$ is chosen to have the form
$$𝒞(\sigma )=\sigma F(u+v)+\frac{1}{\sigma }F(uv),$$
(6.18)
and $`\sigma _0`$ is defined by
$$\sigma =\sigma _0\sqrt{\frac{1+\dot{z}}{1\dot{z}}},$$
(6.19)
then (6.13) becomes a total time derivative of
$$(\dot{z}𝒞𝒰)/\sqrt{1\dot{z}^2}.$$
Hence,
$$𝒫=P+\left[\sigma _0F(u+v)\frac{1}{\sigma _0}F(uv)\right]_{x=z}$$
(6.20)
is exactly conserved. In other words, the contribution to the total momentum of a defect with parameter $`\sigma `$ is precisely the same as it would be if the defect were at rest with parameter $`\sigma _0`$, except that the fields $`u`$ and $`v`$ are evaluated at its actual location $`x=z(t)`$.
Since time translation invariance is unbroken energy conservation is guaranteed with the precise expression given below (6.23). Checking the energy explicitly by calculating the time derivative of the bulk contributions $`E=E(u)+E(v)`$ gives
$$\dot{E}=\left[u_tu_xv_tv_x+\frac{\dot{z}}{2}\left(u_t^2+u_x^2v_t^2v_x^2+2V(u)2W(v)\right)\right]_{x=z}.$$
(6.21)
Using the defect conditions provides a simplification of the right hand side to
$$\frac{d}{dt}\left(\sigma _0F(u+v)+\frac{1}{\sigma _0}F(uv)\right),$$
(6.22)
implying that
$$=E+\left[\sigma _0F(u+v)+\frac{1}{\sigma _0}F(uv)\right]_{x=z}E+𝒞(\sigma _0),$$
(6.23)
is conserved. This is a little surprising since a moving defect might be expected to have some kinetic energy associated with it, which is exchangeable with the fields. However, as was the case with the momentum (6.20), the energy (6.23) is just the energy of the defect system as if it were at rest with parameter $`\sigma _0`$.
In this notation, the potential part of the boundary term is expressed by
$$=\sigma _0(1+\dot{z})F(u+v)+\frac{1}{\sigma _0}(1\dot{z})F(uv),$$
(6.24)
where the fields are evaluated at $`x=z(t)`$.
It has already been remarked that the stationary defects exert no forces on each other and therefore it is expected they might move freely with zero acceleration; besides, from the above discussion, they do not exchange kinetic energy or momentum with the fields. To check this is really consistent, consider the equation of motion for the defect position obtained by varying $`z`$. It turns out to be
$$\left[(u)(v)\right]_{x=z}+\frac{𝐁}{z}\frac{d}{dt}\frac{𝐁}{\dot{z}}=\frac{m_D\ddot{z}}{\sqrt{1\dot{z}^2}}.$$
(6.25)
With the choices made above in (6.7), (6.24), and using the conditions imposed on the fields at the defect by (6.8) and (6.10), the equation of motion (6.25) reduces to
$$m_D\ddot{z}=0.$$
Thus, it is certainly correct that the defect can move with any chosen constant speed and explains why it was not necessary for the defect to exchange kinetic energy with the fields on either side of the defect. If there are several defects, the whole discussion applies independently to each of them. However, because the speed of each defect is arbitrary, defects will pass each other and it becomes necessary to investigate their scattering. If the speed of the defect is itself parameterized using a rapidity variable $`\chi `$ then (6.19) reads
$$\sigma =\mathrm{exp}(\eta )=\mathrm{exp}(\chi \eta _0),\dot{z}=\mathrm{tanh}\chi .$$
(6.26)
Suppose there are two defects with parameters $`\sigma _1`$ and $`\sigma _2`$, and rapidities $`\chi _1>\chi _2`$, initially separating bulk regions containing fields $`u_1,u_2`$ and $`u_2,u_3`$, respectively. A soliton traversing this pair from left to right will be delayed by the first with parameter $`\eta _1`$ and then by the second with parameter $`\eta _2`$. On the other hand, some time later the two defects will have interchanged and a similar soliton would then first encounter the defect with parameter $`\eta _2`$ and then the defect with parameters $`\eta _1`$. In either case, the overall delay will be the same, and this is guaranteed by the Bäckland character of the defect conditions, and in fact represents a statement of Bianchi’s Theorem of Permutability.
The defect conditions (6.9) become very simple with these choices, namely,
$$u_x=\frac{𝒞(\sigma _0)}{u}+v_t,v_x=\frac{𝒞(\sigma _0)}{v}+u_t.$$
(6.27)
This implies straightforwardly, at least for free fields on either side of the defect, that the defect is purely transmitting with the same transmission factor as it would have had the defect been at rest (with parameter $`\eta _0=\eta +\chi `$). A more involved calculation demonstrates that the same is true for the expressions for sine-Gordon soliton delays calculated in section (2).
Returning to the discussion of transmission matrices given in section (3) one might argue as follows: knowing that the transmission matrices depend on the rapidity of the scattering soliton only via $`\theta \eta `$ when the defect is stationary, with parameter $`\eta `$, suggests that they will be expected (according to (6.26)) to depend on $`\theta \eta \chi =\theta \eta _0`$ when the defect has parameter $`\eta `$ but moves with rapidity $`\chi `$. This makes perfect sense, since boosting the scattering soliton by $`\chi `$ would bring the defect with parameter $`\eta `$ relatively to rest. The other parameter $`\nu `$ occurring in (3.55) might then depend upon $`\eta `$, but not upon either $`\theta `$ or $`\chi `$.
In the quantum regime it will be necessary to seek an S-matrix for the defects themselves, call it $`{}_{}{}^{\mathrm{e}}U_{\alpha \beta }^{\gamma \delta }`$, labelled by the initial and final defect topological charges ($`\alpha ,\beta `$) and ($`\gamma ,\delta `$), respectively, satisfying $`\alpha +\beta =\gamma +\delta `$, and compatible with the transmission factors. The defects will be labelled by their parameters and their rapidities so it is first necessary to decide how their scattering matrix elements might depend on these. Suppose their parameters if stationary were $`\eta _{10}=\eta _1+\chi _1,\eta _{20}=\eta _2+\chi _2`$, then the soliton transmission matrix for each defect depends, as argued above, respectively, on $`\theta \eta _{10}`$ and $`\theta \eta _{20}`$. On the other hand, it would be reasonable to suppose the scattering matrix for the defects depends on the difference of their rapidities $`\chi _{12}=\chi _1\chi _2`$, since whether scattering takes place at all depends on the sign of the rapidity difference and not on the relative magnitudes of defect parameters. Besides, a Lorentz transformation effectively shifts $`\theta ,\eta _{10},\eta _{20},\chi _1`$ and $`\chi _2`$, leaving $`\eta _1,\eta _2`$ unchanged.
Using the topological charge labelling, defect-defect scattering will be compatible with the already determined transmission factors provided
$$^\mathrm{e}T_{a\alpha }^{b\gamma }(\theta _1,\eta _1)^\mathrm{e}T_{b\beta }^{c\delta }(\theta _2,\eta _2)^\mathrm{e}U_{\gamma \delta }^{ϵ\rho }(\chi _{12},\eta _1,\eta _2)=^\mathrm{e}U_{\alpha \beta }^{\gamma \delta }(\chi _{12},\eta _1,\eta _2)^\mathrm{e}T_{a\gamma }^{bϵ}(\theta _2,\eta _2)^\mathrm{e}T_{b\delta }^{c\rho }(\theta _1,\eta _1),$$
(6.28)
where $`\theta _1=\theta \eta _{10},\theta _2=\theta \eta _{20}`$. Clearly, $`{}_{}{}^{\mathrm{e}}U`$ is infinite dimensional and, as a consequence of the associativity of (6.28), should itself satisfy a set of Yang-Baxter equations. Since $`\chi _{12}=\eta _{10}\eta _{20}+(\eta _2\eta _1)`$, the dependence of $`{}_{}{}^{\mathrm{e}}U`$ on the defect parameters could be entirely via the rapidity difference but that cannot be assumed initially.
Given the conservation of topological charge, it is convenient to define the elements of $`{}_{}{}^{\mathrm{e}}U`$ as follows,
$$^\mathrm{e}U_{\alpha \beta }^{\gamma \delta }=\underset{\omega }{}\sqrt{\nu _1/\nu _2}^{(\alpha +\beta )/2}\sqrt{\nu _1\nu _2}^{\omega /2}Q^{\omega (\alpha +\beta )/2}A_{\alpha \beta }^\omega \delta _\alpha ^{\gamma +\omega }\delta _\beta ^{\delta \omega },$$
(6.29)
and to make use of the symmetrical version of the transmission matrix given in (3.55). This expression would be simpler if $`\nu _1=\nu _2`$, but this would rule out the possibility of the phase $`\nu `$ having any dependence on the defect parameter. On the other hand, using the alternative representation a basis can always be chosen so that the parameters $`\nu _1,\nu _2`$ are removed.
With this in mind, all the dependence on $`\nu _1`$ or $`\nu _2`$ drops out of the triangle relation and (6.28) reduces to a pair of equations ($`a=c=+;a=+,c=`$) which may be summarized by
$`A_{\alpha +2\beta }^\omega `$ $`=`$ $`A_{\alpha \beta +2}^\omega `$
$`Q^\omega p_1A_{\alpha \beta }^\omega +Q^{\omega 2}p_2A_{\alpha \beta }^{\omega +2}`$ $`=`$ $`Q^\omega p_2A_{\alpha \beta +2}^\omega +Q^{\omega +2}p_1A_{\alpha +2\beta }^{\omega +2},`$ (6.30)
where,
$$p_k=e^{\gamma \eta _{k0}},k=0,1.$$
The other pair ($`a=c=;a=,c=+`$) merely provides the same equations reorganized slightly.
Up to now a neat way to write a solution to this has not been found although clearly the solution depends on the ratio $`p_1/p_2=e^{\gamma (\eta _{10}\eta _{20})}`$. Since unitarity and crossing will also impose constraints, the solution to (6) by itself is not expected to be unique.
The first of (6) states that $`A_{\alpha \beta }^\omega `$ depends only on $`\alpha +\beta `$. Setting $`A_{\alpha \beta }^\omega =B_{\alpha +\beta }^\omega `$, the second of (6) becomes
$$\rho ^2Q^\omega B_\kappa ^\omega +Q^{\omega 2}B_\kappa ^{\omega +2}=Q^\omega B_{\kappa +2}^\omega +\rho ^2Q^{\omega +2}B_{\kappa +2}^{\omega +2},p_1/p_2=\rho ^2.$$
(6.31)
Defining the generating function $`B(y,z)=_{\omega ,\kappa }z^\omega y^\kappa B_\kappa ^\omega `$, the set of equations (6.31) is equivalent to the functional relation
$$B(y,Q^2z)=y^2\left(\frac{1+\rho ^2Q^2z^2}{\rho ^2+Q^2z^2}\right)B(y,z).$$
(6.32)
The latter can be simplified, formally separating the dependence on $`y`$ and $`z`$, by setting
$$B(y,z)=e^{(i\beta ^2/4\pi ^2)\mathrm{ln}y\mathrm{ln}z}a(y)b(z),$$
(6.33)
where $`a(y)`$ is arbitrary, and
$$b(Q^2z)=\left(\frac{1+\rho ^2Q^2z^2}{\rho ^2+Q^2z^2}\right)b(z).$$
(6.34)
The function $`b(z)`$ is assumed to be an analytic function of $`z`$ in some domain excluding the origin and infinity. Since $`|Q|=1`$, if $`z`$ lies on the unit circle the multiplier in (6.34) also lies on the unit circle, and therefore
$$|b(Q^2z)|=|b(z)|.$$
(6.35)
If $`Q^2`$ is a root of unity, iterating (6.34) leads eventually to a contradiction unless $`b(z)`$ is either zero, or diverges for those values of $`Q`$. On the other hand, if $`Q^2`$ is not a root of unity there is no such difficulty. One way to solve this problem might be to multiply $`b(z)`$ by a suitable function of $`Q^2`$ that possesses a set of zeroes at the roots of unity. One example might be the Dedekind eta-function with an appropriate argument, i.e. multiply by a power of
$$\eta (\tau )=e^{i\pi \tau /12}\underset{n=1}{\overset{\mathrm{}}{}}\left(1e^{2i\pi n\tau }\right),\tau =\frac{\gamma +1}{2}=\frac{4\pi }{\beta ^2},$$
(6.36)
and then find a generic solution to (6.34). However, the eta-function is defined for $`\mathrm{Im}\tau >0`$ and a small positive imaginary part would have to be given to $`\gamma `$. \[It is interesting to note that the transformation $`\beta 4\pi i/\beta `$ implements the modular transformation $`\tau 1/\tau `$.\] The eta-function embodies zeroes at the roots of unity but (6.34) gives no hint as to what power of $`\eta (\tau )`$ might be appropriate in these circumstances. A particular solution to (6.34) is given by
$$b(z,\chi _{12})=\underset{k=0}{\overset{\mathrm{}}{}}g_k(\eta _{10}\eta _{20}+\zeta )g_k(\eta _{10}\eta _{20}\zeta ),z=e^{\gamma \zeta /2},$$
(6.37)
where
$$g_k(\lambda )=\frac{\mathrm{\Gamma }\left((k+1/2)\tau +1/4+i\gamma \lambda /4\pi \right)\mathrm{\Gamma }\left((k+1/2)\tau +3/4+i\gamma \lambda /4\pi \right)}{\mathrm{\Gamma }\left((k+1/2)\tau +1/4i\gamma \lambda /4\pi \right)\mathrm{\Gamma }\left((k+1/2)\tau +3/4i\gamma \lambda /4\pi \right)}.$$
(6.38)
This is straightforward to check using standard properties of the gamma-functions. Then, any other solution to (6.34) may only differ from this by a constant factor, at least as far as the $`z`$-dependence is concerned. When $`z`$ lies on the unit circle it is easy to check that the solution given by (6.37) and (6.38) has unit modulus.
An alternative approach would be to label the defects by pairs of integers, as mentioned in section (3), and to denote their scattering matrices by
$$^\mathrm{e}U_{(a,b)(b,c)}^{(a,d)(d,c)}(\chi _{12},\eta _1,\eta _2)^\mathrm{e}U_{abc}^d(\chi _{12}).$$
(6.39)
Clearly, only the middle label changes in the scattering process and the second notation makes use of that fact.
This labelling has an advantage when visualising a two-defect scattering process or the Yang-Baxter equations. In effect, the scattering can be represented by a diagram resembling a ‘quark’ diagram, in which each incoming defect is represented by a pair of lines labelled by the appropriate integers, $`a,b,c,d`$, with the incoming defects sharing $`b`$ and the outgoing defects sharing $`d`$. Then the Yang-Baxter equation asserts the equality of a pair of diagrams each containing a single closed loop.
In this representation, the triangle relation for $`a=c=+`$ gives rise to two equations, which may be written:
$$^\mathrm{e}U_{abc}^d(\chi _{12})=^\mathrm{e}U_{a+1b+1c+1}^{d+1}(\chi _{12})=^\mathrm{e}U_{a1b+1c1}^{d+1}(\chi _{12}).$$
(6.40)
From these, it is clear that
$${}_{}{}^{\mathrm{e}}U_{abc}^{d}(\chi _{12})=^\mathrm{e}U_{\mathrm{0\; 0}ca}^{db}(\chi _{12}),$$
where, in terms of the topological charges, $`ca=\alpha +\beta `$ and $`bd=\alpha \gamma `$. Therefore, the alternative labelling does not appear to provide any more generality than the labelling in terms of topological charges.
Apart from (6.34), there are other conditions arising from the requirements of unitarity and crossing, or the possibility that defects of opposite charge may annihilate virtually to the vacuum. In terms of $`{}_{}{}^{\mathrm{e}}U`$ the unitarity condition (for $`\chi `$ real) is,
$$^\mathrm{e}U_{\alpha \beta }^{\gamma \delta }(\chi )^\mathrm{e}\overline{U}_{\delta \gamma }^{ϵ\rho }(\chi )=\delta _\alpha ^\rho \delta _\beta ^ϵ,$$
(6.41)
and this translates in terms of $`B_\kappa ^\omega `$ into the collection of relations,
$`{\displaystyle \underset{\omega }{}}B_\kappa ^\omega \overline{B}_\kappa ^{\omega +ϵ}=\delta _0^ϵ,`$ (6.42)
which ought to hold for each choice of $`ϵ`$ and $`\kappa `$. This may be rewritten in terms of generating functions setting $`\overline{B}(y,z)=_{\omega \kappa }z^\omega y^\kappa \overline{B}_\kappa ^\omega `$. If both $`y`$ and $`z`$ lie on unit circles, $`\overline{B}(y,z)`$ is the complex conjugate of $`B(y,z)`$. Then, using the separation property of $`B(y,z)`$ given in (6.33), the unitarity condition itself separates into two types of term
$$\frac{dz}{2\pi iz}z^ϵb(z)\overline{b}(z)=0,\text{for}ϵ0,$$
(6.43)
and
$$\frac{dz}{2\pi iz}b(z)\overline{b}(z)\frac{dy}{2\pi iy}y^\kappa a(y)\frac{dy^{}}{2\pi iy^{}}y^\kappa \overline{a}(y^{})=1,\text{for each}\kappa .$$
(6.44)
In all cases it has been assumed there is a singularity-free region enclosing the origin and containing the contours of integration. Equation (6.43) follows automatically on choosing the $`z`$ contour to lie on the unit circle since it has been noted already that $`b(z)`$ has unit modulus there. The other part of the condition is more problematical since it appears to require that each coefficient in a Laurent expansion of $`a(y)`$ is equal to the inverse of the corresponding coefficient in a similar expansion of $`\overline{a}(y)`$, at least up to an overall, term-independent constant.
It is tempting to suppose the crossing property should take the form
$$U_{\alpha \beta }^{\gamma \delta }(i\pi \chi _{12},\eta _1,\eta _2)=U_{\alpha \gamma }^{\beta \delta }(\chi _{12},\eta _1,\eta _2),$$
(6.45)
which, if $`\nu _1=\nu _2=\nu `$ translates in terms of $`B_\kappa ^\omega `$ to,
$$B_\kappa ^\omega (i\pi \chi _{12})=\nu ^{(\omega \kappa )/2}B_\omega ^\kappa (\chi _{12}).$$
(6.46)
Because the upper and lower indices on $`B_\kappa ^\omega `$ are interchanged it will be possible to use the crossing relation to relate the two functions $`a(y)`$ and $`b(z)`$. Rewriting (6.46) in terms of the generating functions leads to the relation,
$$\frac{a(y,i\pi \chi _{12})}{b(y/\sqrt{\nu },\chi _{12})}e^{(i\beta ^2/8\pi ^2)\mathrm{ln}\nu \mathrm{ln}y}=\frac{a(\sqrt{\nu }z,\chi _{12})}{b(z,i\pi \chi _{12})}e^{(i\beta ^2/8\pi ^2)\mathrm{ln}\nu \mathrm{ln}\sqrt{\nu }z},$$
(6.47)
from which it is clear both sides must be independent of $`y`$ and $`z`$, and separately equal to a crossing symmetric function $`c(\chi _{12})=c(i\pi \chi _{12})`$. Then, either side of (6.47) implies an expression for the unkown function $`a(y)`$ in terms of $`b(y)`$, itself given by (6.37). Explicitly, the expression for $`a(y)`$ is,
$$a(y,\chi _{12})=c(\chi _{12})b(y/\sqrt{\nu },i\pi \chi _{12})e^{(i\beta ^2/8\pi ^2)\mathrm{ln}\nu \mathrm{ln}y}.$$
(6.48)
Clearly, the expression (6.48) simplifies significantly if $`\nu =1`$.
Alternatively, the analogue of (3.34) representing the possibility for a defect to annihilate virtually with a defect of opposite topological charge would be
$$c_{\beta \beta }^0\delta _\alpha ^\rho =\underset{ϵ,\gamma }{\overset{\mathrm{e}}{}}U_{\alpha \beta }^{ϵ\gamma }(\chi i\pi /2)^\mathrm{e}U_{\gamma \beta }^{ϵ\rho }(\chi +i\pi /2)c_{ϵϵ}^0,$$
(6.49)
and in terms of $`B_\kappa ^\omega `$ this becomes
$$c_{\beta \beta }^0=\underset{ϵ}{}Q^{ϵ\alpha }\nu ^{(ϵ+\alpha )/2}B_\alpha ^ϵ(\chi i\pi /2)B_ϵ^\alpha (\chi +i\pi /2)c_{\alpha \beta ϵϵ+\beta \alpha }^0.$$
(6.50)
However, without detailed knowledge of the couplings it is not easy to extract information from this relationship.
## 7 Discussion
The jump-defect possesses a variety of interesting classical properties stemming from the fact that the integrable defect conditions are a ‘frozen’ Bäcklund transformation. The purpose of this article has been to explore the corresponding quantum sine-Gordon field theory and investigate the extent to which the classical picture extends to the quantum domain. One of the striking features of the classical scattering of a soliton with a jump-defect was the possibility of the soliton being ‘eaten’ by the defect, and another was the possibility of the soliton converting to an anti-soliton. The fact a soliton can disappear and be replaced by an ‘excited’ defect, carrying exactly the energy and momentum of the soliton, emerges in the quantum field theory as a resonant soliton-defect bound state, with a finite, coupling-dependent width. Scattering with the excited state is not expected to be unitary (since the excited defect is no longer a possible asymptotic state) and that turns out to be precisely the case. Curiously, the defect treats solitons and anti-solitons differently. The algebraic analysis of the triangle relations makes this clear. Morever, an argument presented in section (3), based on a functional integral representation of the transmission matrix, demonstrates how even this fact follows from the Lagrangian starting point (2.2). Another surprising feature appears when one considers the breather transmission factors. Each of these may be defined via a bootstrap procedure. However, the transmission factor for the lightest breather is predicted to be completely independent of the coupling constant. This fact, if correct, should be verifiable in perturbation theory by calculating the quantum corrections to the classical result given in section (2). On the other hand, if the prediction is false, it would indicate the necessity of additional CDD factors in the derived transmission matrix presented in (3.55), and a detailed perturbative calculation may provide a pointer towards the precise form of any such factors. Typically, the other breather transmission factors have poles at complex rapidity whose interpretation is obscure owing to the lack of classical configurations that might correspond to them. Possibly they indicate a collection of new excited states of the even-charged defect whose details will depend upon the value of the bulk coupling. Or, they may be artifacts of the bootstrap procedure. In either case, a detailed explanation is needed.
There has been work previously on models containing unstable bound states , including interesting ideas concerning their consequences in the context of impurities . However, it appears the present work may provide the first example of unstable states appearing naturally within the sine-Gordon model itself. For the reasons stated, it is highly desirable to find a physical situation where the sine-Gordon Bäcklund transformation is part of the formulation. If it exists a physical system of this kind would offer the opportunity to control solitons and have possible consequences for communications making use of solitons.
There is no reason why defects should not be able to move with constant speed, and therefore to scatter. Classically, the delays, or changes of character, suffered by a soliton scattering with two moving defects is entirely consistent once the defects have changed places. In fact the consistency is an expression of Bianchi’s identity concerning the permutability of a pair of Bäcklund transformations. The classical picture of a moving defect is analyzed in detail and it has been verified that with a natural choice of defect condition, the jump-defect moves at constant speed. Moreover, multiple defects have no long range interactions with each other. For this reason it is plausible that their interaction in the quantum domain is local and factorisable, implying that there should be a triangle compatibility relation (6.28), involving the scattering matrix for a pair of defects and the already derived transmission matrices for the solitons. The triangle relation is supplemented by unitarity and crossing relations and a fuller discussion will be given in a subsequent paper. In this regard, the purpose of this article has been to develop the framework within which such questions may be asked and to give an indication of how they may be tackled. Further investigation will be needed to elaborate all the details.
Since the sine-Gordon model allows both integrable boundary conditions and integrable jump-defects one might also ask how defects behave when they encounter a boundary. To some extent, the question should be answerable algebraically by finding solutions to the associated reflection Yang-Baxter equations , once the defect scattering matrix is completely determined - although such solutions would need to be supplemented by other information. One might speculate that placing a stationary defect close to a boundary could modify the boundary condition. The jump-defects provide an environment where such questions can be asked, and answered, up to a point. In the $`\mathrm{sinh}`$-Gordon model all the integrable boundary conditions and reflection factors are known . However, the scalar particle has a $`\beta `$-independent transmission factor that would appear to be quite unable to account for the intricate $`\beta `$-dependence of the reflection factors for general boundary conditions, starting from, for example, the reflection factor associated with a Neumann boundary condition. That is, of course, presuming the results of section (5) are entirely unmodified by CDD factors. Otherwise, it might be possible to invert the argument and use known reflection factors to inform the discussion of transmission factors.
A more interesting question for the future will be to decide whether the jump-defects themselves might be described by a local quantum field theory coupled to (or, even a part of) the sine-Gordon model. This was hinted at in the introduction where it was pointed out that a pair of stationary, classical, jump-defects behave like a soliton, at least in so far as mimicking the delay a soliton would experience overtaking another soliton. It might be fruitful to regard solitons as bound pairs of defects. In which case, whatever the theory of defects is, it should contain the sine-Gordon model, perhaps as a suitable limit. One might also wonder about the dual relationship between the sine-Gordon model and the massive Thirring model, formulated by Skyrme, Coleman and Mandelstam . If the jump-defects have a role to play within the sine-Gordon model, then they should also be describable within the massive Thirring model where the basic fields are fermions equivalent to the solitons of the sine-Gordon model. One might imagine a hierarchy in which the defects bind to form solitons, and the solitons bind to form breathers. This aspect needs further investigation.
There are other integrable models with classical soliton behaviour and it will be interesting to see if the jump-defect idea may be incorporated naturally within those models. Some results are already known for the affine Toda field theories, especially those of $`a_n`$ type . However, the results known so far are restricted to classical integrability, and placing the jump-defects in a quantum context for these examples has yet to be accomplished.
Acknowledgements CZ thanks JSPS for a Fellowship and EC is indebted to members of the Ecole Normale Superieure de Lyon, the University of Bologna and the Yukawa Institute for Theoretical Physics, and especially Jean Michel Maillet, Francesco Ravanini and Ryu Sasaki, for their hospitality. PB thanks Patrick Dorey for discussions and EC wishes to thank Davide Fioravanti for comments. The work has also been supported by the Leverhulme Trust and EUCLID - a European Commission RTN Network (contract HPRN-CT-2002-00325). EC would like to thank the organisers of the 2nd EUCLID Annual Conference (September 2004, Sozopol, Bulgaria) and PB thanks the organisers of the 3rd EUCLID Spring School (May 2005, Trieste, Italy) for the opportunity to present some of the ideas discussed in this article.
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# Combinatorial Interpretations of the 𝑞-Faulhaber and 𝑞-Salié Coefficients
## 1. Introduction
In the early seventeenth century, Johann Faulhaber (see also ) considered the sums of powers $`S_{m,n}=_{k=1}^nk^m`$ and provided formulas for the coefficients $`f_{m,k}`$ ($`0m8`$) in
$$S_{2m+1,n}=\frac{1}{2}\underset{k=1}{\overset{m}{}}f_{m,k}\left(n(n+1)\right)^{k+1},$$
(1)
In 1989, Ira Gessel and Xavier Viennot studied the alternating sum $`T_{2m,n}=_{k=1}^n(1)^{nk}k^{2m}`$ and showed that there exist integers $`s_{m,k}`$ such that
$$T_{2m,n}=\frac{1}{2}\underset{k=1}{\overset{m}{}}s_{m,k}(n(n+1))^k.$$
(2)
In particular, they proved that the Faulhaber coefficients $`f_{m,k}`$ and the Salié coefficients $`s_{m,k}`$ count certain families of non-intersecting lattice paths.
Recently, two of the authors , continuing work of Michael Schlosser , Sven Ole Warnaar and Kristina Garrett and Kristen Hummel , have found $`q`$-analogues of (1) and (2). More precisely, setting $`[k]=\frac{1q^k}{1q}`$, $`[k]!=_{i=1}^k[k]`$, and
$`S_{m,n}(q)`$ $`={\displaystyle \underset{k=1}{\overset{n}{}}}{\displaystyle \frac{[2k]}{[2]}}[k]^{m1}q^{\frac{m+1}{2}(nk)},`$ (3)
$`T_{m,n}(q)`$ $`={\displaystyle \underset{k=1}{\overset{n}{}}}(1)^{nk}[k]^mq^{\frac{m}{2}(nk)},`$ (4)
for $`m,n`$, they proved the following results:
###### Theorem 1.1.
There exist polynomials $`P_{m,k}`$, $`Q_{m,k}`$, $`G_{m,k}`$ and $`H_{m,k}`$ in $`[q]`$ such that
$`S_{2m+1,n}(q)`$ $`={\displaystyle \underset{k=0}{\overset{m}{}}}(q^n)^{mk}{\displaystyle \frac{[k]!}{[m+1]!}}P_{m,mk}(q){\displaystyle \frac{([n][n+1])^{k+1}}{[2]}},`$ (5)
$`S_{2m,n}(q)`$ $`=(1q^{n+\frac{1}{2}}){\displaystyle \underset{k=0}{\overset{m}{}}}(q^n)^{mk}{\displaystyle \frac{(1q^{\frac{1}{2}})^{mk}Q_{m,mk}(q^{\frac{1}{2}})}{_{i=0}^{mk}(1q^{mi+\frac{1}{2}})}}{\displaystyle \frac{([n][n+1])^k}{[2]}},`$ (6)
$`T_{2m,n}(q)`$ $`={\displaystyle \underset{k=1}{\overset{m}{}}}(q^n)^{mk}{\displaystyle \frac{G_{m,mk}(q)}{_{i=0}^{mk}(1+q^{mi})}}([n][n+1])^k,`$ (7)
and
$`T_{2m1,n}(q)`$ $`=(1)^{m+n}H_{m,m1}(q^{\frac{1}{2}}){\displaystyle \frac{q^{(m\frac{1}{2})n}}{(1+q^{\frac{1}{2}})^m_{i=0}^{m1}(1+q^{mi\frac{1}{2}})}}`$
$`+{\displaystyle \frac{1q^{n+\frac{1}{2}}}{1q^{\frac{1}{2}}}}{\displaystyle \underset{k=1}{\overset{m}{}}}(q^n)^{mk}{\displaystyle \frac{H_{m,mk}(q^{\frac{1}{2}})([n][n+1])^{k1}}{(1+q^{\frac{1}{2}})^{mk+1}_{i=0}^{mk}(1+q^{mi\frac{1}{2}})}}.`$ (8)
Comparing with (3) and (4), we have
$`f_{m,k}`$ $`=(1)^{mk}{\displaystyle \frac{k!}{(m+1)!}}P_{m,mk}(1)`$
and
$`s_{m,k}`$ $`=(1)^{mk}2^{km}G_{m,mk}(1),`$
but the numbers corresponding to $`Q_{m,k}(1)`$ and $`H_{m,k}(1)`$ do not seem to be studied in the literature. The first values of $`P_{m,k}`$, $`Q_{m,k}`$, $`G_{m,k}`$ and $`H_{m,k}`$ are given in Tables 14, respectively.
Recall that a polynomial $`f(x)=a_0+a_1x+\mathrm{}+a_nx^n`$ of degree $`n`$ has *symmetric coefficients* if $`a_i=a_{ni}`$ for $`0in`$. The tables above suggest that the coefficients of the polynomials $`P_{m,k}`$, $`Q_{m,k}`$, $`G_{m,k}`$ and $`H_{m,k}`$ are nonnegative and symmetric. The aim of this paper is to prove this fact by showing that the coefficients count certain families of non-intersecting lattice paths.
## 2. Inverses of matrices
Recall that the $`n`$-*th complete homogeneous functions in $`r`$ variables* $`x_1,x_2,\mathrm{},x_r`$ has the following generating function:
$$\underset{n0}{}h_n(x_1,\mathrm{},x_r)t^n=\frac{1}{(1x_1t)(1x_2t)\mathrm{}(1x_rt)}.$$
For $`r,s0`$, let $`h_n(\{1\}^r,\{q\}^s)`$ denote the $`n`$-th complete homogeneous functions in $`r+s`$ variables, of which $`r`$ are specialized to $`1`$ and the others to $`q`$, i.e.,
$$\underset{n0}{}h_n(\{1\}^r,\{q\}^s)z^n=\frac{1}{(1z)^r(1qz)^s}.$$
(9)
By convention, $`h_n(\{1\}^r,\{q\}^s)=0`$ if $`r<0`$ or $`s<0`$. For convenience, we also write $`h_n(\{1,q\}^r)`$ instead of $`h_n(\{1\}^r,\{q\}^r)`$.
We first prove the following result.
###### Lemma 2.1.
Let $`a`$ and $`b`$ be non-negative integers, then
$$\begin{array}{c}\underset{m0}{}\underset{k0}{}h_{m2k}(\{1\}^{k+a},\{q\}^{k+b})\left(\frac{q^l}{[l]^2}\right)^kz^m=\frac{[l]^2}{[2l]}\{\begin{array}{cc}\frac{[l+1]}{[l][l+1]z}\frac{q[l1]}{[l]q[l1]z}\hfill & \text{for }a=1\text{}b=1\text{,}\hfill \\ \frac{1}{[l][l+1]z}+\frac{q^l}{[l]q[l1]z}\hfill & \text{for }a=1\text{}b=0\text{}\hfill \\ \frac{q^l}{[l][l+1]z}+\frac{1}{[l]q[l1]z}\hfill & \text{for }a=0\text{}b=1\text{.}\hfill \end{array}\hfill \end{array}$$
###### Proof.
Using the definition (9) of the complete homogeneous functions we have
$$\begin{array}{cc}& \underset{m0}{}\underset{k0}{}h_{m2k}(\{1\}^{k+a},\{q\}^{k+b})x^kz^m\hfill \\ & =\underset{k0}{}\frac{x^kz^{2k}}{(1z)^{k+a}(1qz)^{k+b}}\hfill \\ & =\frac{1}{(1z)^{a1}(1qz)^{b1}}\frac{1}{(1z)(1qz)xz^2}.\hfill \end{array}$$
Setting $`x=\frac{q^l}{[l]^2}`$ a little calculation shows that the denominator of the second fraction factorizes:
$$\frac{1}{(1z)(1qz)xz^2}=\frac{1}{\left([l]qz[l1]\right)\left([l]z[l+1]\right)}.$$
The result then follows from the standard partial fraction decomposition. ∎
Let $`X_n=\frac{[n][n+1]}{q^n}`$. The following lemma might be interesting per se. When $`q=1`$ it reduces to simple applications of the binomial theorem.
###### Lemma 2.2.
For $`k,m1`$, set
$`c_{k,m}(q)`$ $`:=h_{2mk}(\{1,q^2\}^{km+1})+qh_{2mk1}(\{1,q^2\}^{km+1}),`$
$`g_{k,m}(q)`$ $`:=h_{2mk}(\{1\}^{km+1},\{q\}^{km})+h_{2mk}(\{1\}^{km},\{q\}^{km+1}),`$
$`d_{k,m}(q)`$ $`:=g_{k,m}(q^2)+qg_{k1,m1}(q^2).`$
For $`m,l1`$, we have
$`X_l^{m+1}X_{l1}^{m+1}`$ $`={\displaystyle \underset{k}{}}h_{m2k}(\{1,q\}^{k+1})[2l][l]^{2(mk)}q^{l(mk+1)},`$ (10)
$`{\displaystyle \frac{1q^{l+\frac{1}{2}}}{(1q^{\frac{1}{2}})q^{\frac{l}{2}}}}X_l^m{\displaystyle \frac{1q^{l\frac{1}{2}}}{(1q^{\frac{1}{2}})q^{\frac{l1}{2}}}}X_{l1}^m`$ $`={\displaystyle \underset{k}{}}c_{m,mk}(q^{\frac{1}{2}})[2l][l]^{2(mk\frac{1}{2})}q^{l(mk+\frac{1}{2})},`$ (11)
$`X_l^m+X_{l1}^m`$ $`={\displaystyle \underset{k}{}}g_{m,mk}(q)[l]^{2(mk)}q^{l(mk)},`$ (12)
$`{\displaystyle \frac{1q^{l+\frac{1}{2}}}{(1q^{\frac{1}{2}})q^{\frac{l}{2}}}}X_l^{m1}+{\displaystyle \frac{1q^{l\frac{1}{2}}}{(1q^{\frac{1}{2}})q^{\frac{l1}{2}}}}X_{l1}^{m1}`$ $`={\displaystyle \underset{k}{}}d_{m,mk}(q^{\frac{1}{2}})[l]^{2(mk\frac{1}{2})}q^{l(mk\frac{1}{2})}.`$ (13)
###### Proof.
The proof rests on the previous lemma.
* Applying Lemma 2.1 with $`a=1`$ and $`b=1`$ yields that the coefficient of $`z^m`$ in $`_kh_{m2k}(\{1,q\}^{k+1})q^{lk}[l]^{2k}`$ is
$$\frac{[l]}{[2l]}\left([l+1]\left(\frac{[l+1]}{[l]}\right)^mq[l1]\left(\frac{q[l1]}{[l]}\right)^m\right).$$
Multiplying this expression with $`[2l]\frac{[l]^{2m}}{q^{l(m+1)}}`$ we obtain (10).
* Since $`c_{m,mk}(q^{\frac{1}{2}})=h_{m2k}(\{1,q\}^{k+1})+q^{\frac{1}{2}}h_{m12k}(\{1,q\}^{k+1})`$, Equation (11) follows directly from the previous calculation.
* As $`g_{m,mk}(q)=h_{m2k}(\{1\}^{k+1},\{q\}^k)+h_{m2k}(\{1\}^k,\{q\}^{k+1})`$, applying Lemma 2.1 with $`a=1`$, $`b=0`$ and $`a=0`$, $`b=1`$,
$`{\displaystyle \underset{k}{}}(h_{m2k}(\{1\}^{k+1},\{q\}^k)+h_{m2k}(\{1\}^k,\{q\}^{k+1})))q^{lk}[l]^{2k}`$
$`={\displaystyle \frac{[l]^2}{[2l]}}\left({\displaystyle \frac{1+q^l}{[l][l+1]z}}+{\displaystyle \frac{1+q^l}{[l]q[l1]z}}\right)`$
Multiplying the coefficient of $`z^m`$ of this expression with $`[l]^{2m}q^{lm}`$ we obtain (12).
* Since $`d_{m,mk}(q^{\frac{1}{2}})=g_{m,mk}(q)+q^{\frac{1}{2}}g_{m1,mk1}(q)`$, Equation (13) follows directly from the previous calculation.
The following is the main result of this section. Note that together with Theorems 3.2 and 3.6 it also provides an alternative proof of Theorem 1.1.
###### Theorem 2.3.
The inverses of the lower triangular matrices
$$(h_{2mk}(\{1,q\}^{km+1}))_{0k,mn},(c_{k,m}(q))_{1k,mn},(g_{k,m}(q))_{1k,mn},(d_{k,m}(q))_{1k,mn}$$
are respectively the lower triangular matrices
$`\left((1)^{km}{\displaystyle \frac{[m]!}{[k+1]!}}P_{k,km}(q)\right)_{0k,mn},`$ (14)
$`\left((1)^{km}{\displaystyle \frac{(1q)^{km+1}Q_{k,km}(q)}{_{i=0}^{km}(1q^{2k2i+1})}}\right)_{1k,mn},`$ (15)
$`\left((1)^{km}{\displaystyle \frac{G_{k,km}(q)}{_{i=0}^{km}(1+q^{ki})}}\right)_{1k,mn},`$ (16)
$`\left((1)^{km}{\displaystyle \frac{H_{k,km}(q)}{(1+q)^{km+1}_{i=0}^{km}(1+q^{2k2i1})}}\right)_{1k,mn}.`$ (17)
###### Proof.
* Summing Equation (10) over $`l`$ from $`1`$ to $`n`$ and applying Equation (3), we obtain
$`X_n^{m+1}=[2]{\displaystyle \underset{k=0}{\overset{m/2}{}}}h_{m2k}(\{1,q\}^{k+1})S_{2m2k+1,n}(q)q^{n(mk+1)}.`$ (18)
Plugging (5) in Equation (18), the right-hand side becomes
$$\underset{k=0}{\overset{m/2}{}}\underset{l=0}{\overset{mk}{}}h_{m2k}(\{1,q\}^{k+1})(1)^{mkl}\frac{[l]!}{[mk+1]!}P_{mk,mkl}(q)X_n^{l+1}.$$
Comparing the coefficients of $`X_n^{l+1}`$ we see that $`(h_{2mk}(\{1,q\}^{km+1}))_{0k,mn}`$ and (14) are indeed inverses.
* Summing Equation (11) over $`l`$ from $`1`$ to $`n`$ and applying Equation (3), we obtain
$`{\displaystyle \frac{1q^{n+\frac{1}{2}}}{(1q^{\frac{1}{2}})q^{\frac{n}{2}}}}X_n^m`$ $`=[2]{\displaystyle \underset{k=0}{\overset{m/2}{}}}c_{m,mk}(q^{\frac{1}{2}})S_{2m2k,n}(q)q^{n(mk+\frac{1}{2})}.`$ (19)
Substituting (6) into (19) and dividing both sides by $`\frac{1q^{n+\frac{1}{2}}}{(1q^{\frac{1}{2}})q^{\frac{n}{2}}}`$, we get
$$X_n^m=\underset{k=0}{\overset{m/2}{}}\underset{l=1}{\overset{mk}{}}c_{m,mk}(q^{\frac{1}{2}})(1)^{mkl}\frac{(1q^{\frac{1}{2}})^{mkl}Q_{mk,mkl}(q^{\frac{1}{2}})}{_{i=0}^{mkl}(1q^{mki+\frac{1}{2}})}X_n^l.$$
Comparing the coefficients of $`X_n^l`$, we see that $`(c_{k,m}(q))_{1k,mn}`$ and (15) are indeed inverses.
* Equation (12) may be written as
$`(1)^{nl}X_l^m(1)^{nl+1}X_{l1}^m=(1)^{nl}{\displaystyle \underset{k=0}{\overset{m/2}{}}}g_{m,mk}(q){\displaystyle \frac{(1q^l)^{2m2k}}{(1q)^{2m2k}}}q^{l(mk)}`$ (20)
Summing Equation (20) over $`l`$ from $`1`$ to $`n`$ and applying Equation (4), we obtain
$`X_n^m={\displaystyle \underset{k=0}{\overset{m/2}{}}}g_{m,mk}(q)T_{2m2k,n}(q)q^{n(mk)}.`$ (21)
Substituting (7) into (21), the right-hand side becomes
$`{\displaystyle \underset{k=0}{\overset{m/2}{}}}{\displaystyle \underset{l=1}{\overset{mk}{}}}g_{m,mk}(q)(1)^{mkl}{\displaystyle \frac{G_{mk,mkl}(q)}{_{i=0}^{mkl}(1+q^{mki})}}X_n^l.`$ (22)
Comparing the coefficients of $`X_n^l`$, we see that $`(g_{k,m}(q))_{1k,mn}`$ and (16) are inverse to each other.
* Equation (13) may be written as
$`(1)^{nl}{\displaystyle \frac{1q^{l+\frac{1}{2}}}{(1q^{\frac{1}{2}})q^{\frac{l}{2}}}}X_l^{m1}(1)^{nl+1}{\displaystyle \frac{1q^{l\frac{1}{2}}}{(1q^{\frac{1}{2}})q^{\frac{l1}{2}}}}X_{l1}^{m1}`$
$`=(1)^{nl}{\displaystyle \underset{k}{}}d_{m,mk}(q^{\frac{1}{2}})[l]^{2(mk\frac{1}{2})}q^{l(mk\frac{1}{2})}.`$ (23)
Summing Equation (23) over $`l`$ from $`1`$ to $`n`$ and applying Equation (4), we obtain
$`{\displaystyle \frac{1q^{n+\frac{1}{2}}}{(1q^{\frac{1}{2}})q^{\frac{n}{2}}}}X_n^{m1}={\displaystyle \underset{k}{}}d_{m,mk}(q^{\frac{1}{2}})T_{2m2k1,n}(q)q^{n(mk\frac{1}{2})},m2.`$ (24)
Substituting (1.1) into (24) yields
$`{\displaystyle \frac{1q^{n+\frac{1}{2}}}{(1q^{\frac{1}{2}})q^{\frac{n}{2}}}}\left(X_n^{m1}{\displaystyle \underset{k}{}}{\displaystyle \underset{l=1}{\overset{mk}{}}}{\displaystyle \frac{(1)^{mkl}d_{m,mk}(q^{\frac{1}{2}})H_{mk,mkl}(q^{\frac{1}{2}})X_n^{l1}}{(1+q^{\frac{1}{2}})^{mkl+1}_{i=0}^{mkl}(1+q^{mki\frac{1}{2}})}}\right)`$
$`=(1)^n{\displaystyle \underset{k}{}}{\displaystyle \frac{(1)^{mk}d_{m,mk}(q^{\frac{1}{2}})H_{mk,mk1}(q^{\frac{1}{2}})}{(1+q^{\frac{1}{2}})^{mk}_{i=0}^{mk1}(1+q^{mki\frac{1}{2}})}}.`$ (25)
We now show that the right-hand side of (25) must vanish. Suppose $`0<q<1`$. Denote the left-hand side of (25) by $`L_n`$. If there exists an $`n`$ such that $`L_n=0`$ we are done. Suppose $`L_n0`$ for all $`n1`$, then $`L_n`$ is a rational function in $`t=q^{\frac{n}{2}}`$ and can be written as
$$L_n=t^sf(t)\text{with}t=q^{\frac{n}{2}},$$
where $`s`$ is an integer and $`f(t)`$ a rational function with $`f(0)0`$. Since $`f(q^{\frac{n}{2}})0`$, the right-hand side of (25) implies that
$$f(q^{\frac{n}{2}})f(q^{\frac{n+1}{2}})<0n1.$$
Taking the limit as $`n\mathrm{}`$ we get $`(f(0))^20`$, which is impossible. Hence $`L_n=0`$ and (25) reduces to
$`X_n^{m1}={\displaystyle \underset{k}{}}d_{m,mk}(q^{\frac{1}{2}}){\displaystyle \underset{l=1}{\overset{mk}{}}}{\displaystyle \frac{(1)^{mkl}H_{mk,mkl}(q^{\frac{1}{2}})X_n^{l1}}{(1+q^{\frac{1}{2}})^{mkl+1}_{i=0}^{mkl}(1+q^{mki\frac{1}{2}})}}.`$ (26)
Comparing the coefficients of $`X_n^{l1}`$ on both sides of (26), we see that $`(d_{k,m}(q))_{1k,mn}`$ and (17) are indeed inverses.
The following easily verified result has been given by Gessel and Viennot .
###### Lemma 2.4.
Let $`(A_{ij})_{0i,jm}`$ be an invertible lower triangular matrix, and let $`(B_{ij})=(A_{ij})^1`$. Then for $`0knm`$, we have
$$B_{n,k}=\frac{(1)^{nk}}{A_{k,k}A_{k+1,k+1}\mathrm{}A_{n,n}}\left|A_{k+i+1,k+j}\right|_{0i,jnk1}.$$
Using the above lemma we derive immediately from Theorem 2.3 the following determinant formulas:
$`P_{m,k}(q)`$ $`=\underset{0i,jk1}{det}(h_{mki+2j1}(\{1,q\}^{ij+2})),`$ (27)
$`Q_{m,k}(q)`$ $`=\underset{0i,jk1}{det}(c_{mk+i+1,mk+j}(q)),`$ (28)
$`G_{m,k}(q)`$ $`=\underset{0i,jk1}{det}(g_{mk+i+1,mk+j}(q)),`$ (29)
$`H_{m,k}(q)`$ $`=\underset{0i,jk1}{det}(d_{mk+i+1,mk+j}(q)).`$ (30)
## 3. Combinatorial interpretations
A *lattice path* or *path* $`s_0s_n`$ is a sequence of points $`(s_0,s_1,\mathrm{},s_n)`$ in the plane $`^2`$ such that $`s_is_{i1}=(1,0),(0,1)`$ for all $`i=1,\mathrm{},n`$. Let us assign a weight to each step $`(s_i,s_{i+1})`$ of $`s_0s_n`$. We define the weight $`N(s_0s_n)`$ of the path $`s_0s_n`$ to be the product of the weights of its steps. Let $`s_0=(a,b)`$ and $`s_n=(c,d)`$, if we weight each vertical step with $`x`$-coordinate $`i`$ by $`x_i`$ and all horizontal steps by 1 then
$$N(s_0s_n)=h_{db}(x_a,x_{a+1},\mathrm{},x_c).$$
(31)
Now consider two sequences of lattice points $`𝐮:=(u_1,u_2,\mathrm{},u_n)`$ and $`𝐯:=(v_1,v_2,\mathrm{},v_n)`$ such that for $`i<j`$ and $`k<l`$ any lattice path between $`u_i`$ and $`v_l`$ has a common point with any lattice path between $`u_j`$ and $`v_k`$. Set
$$N(𝐮,𝐯):=N(u_1v_1)\mathrm{}N(u_nv_n),$$
where the sum is over all families of non-intersecting paths $`(u_1v_1,\mathrm{},u_nv_n)`$.
The following remarkable result can be found in Gessel and Viennot . For historical remarks see also Krattenthaler .
###### Theorem 3.1.
\[Lindström-Gessel-Viennot\] We have
$$N(𝐮,𝐯)=\underset{1i,jn}{det}(N(u_jv_i)).$$
We are now ready to exhibit the combinatorial interpretation of the $`q`$-Faulhaber numbers.
###### Theorem 3.2.
Let $`𝐮=(u_0,\mathrm{},u_{k1})`$ and $`𝐯=(v_0,\mathrm{},v_{k1})`$, where $`u_i=(2i,2i)`$ and $`v_i=(2i+3,mki1)`$ for $`0ik1`$.
1. The polynomial $`P_{m,k}(q)`$ is the sum of the weights of $`k`$-non-intersecting paths from $`𝐮`$ to $`𝐯`$, where a vertical step with an even $`x`$-coordinate has weight $`q`$, and all the other steps have weight $`1`$.
2. The polynomial $`Q_{m,k}(q)`$ is the sum of the weights of $`k`$-non-intersecting paths from $`𝐮`$ to $`𝐯`$, where the weight of the individual steps is the same as before with the exception that $`q`$ is replaced with $`q^2`$ and the vertical step starting from any $`u_j`$ has weight $`q^2+q`$ instead of $`q^2`$.
###### Proof.
For (i), by means of (31) we have
$$N(u_jv_i)=h_{mki+2j1}(\{1,q\}^{ij+2}).$$
The result then follows from (27) and Theorem 3.1.
For (ii), assume that $`u_j^{}=(2j+1,2j)`$ and $`u_j^{\prime \prime }=(2j,12j)`$. The first step of a lattice path from $`u_j`$ to $`v_i`$ is either $`u_ju_j^{}`$ or $`u_ju_j^{\prime \prime }`$. As $`N(u_ju_j^{})=1`$, $`N(u_ju_j^{\prime \prime })=q^2+q`$ and $`h_n(x_1,\mathrm{},x_{r1})+x_rh_{n1}(x_1,\mathrm{},x_r)=h_n(x_1,\mathrm{},x_r)`$, we have
$`N(u_jv_i)`$ $`=N(u_ju_j^{})N(u_j^{}v_i)+N(u_ju_j^{\prime \prime })N(u_j^{\prime \prime }v_i)`$
$`=N(u_j^{}v_i)+(q^2+q)N(u_j^{\prime \prime }v_i)`$
$`=h_{mki+2j1}(\{1\}^{ij+2},\{q^2\}^{ij+1})`$
$`+(q^2+q)h_{mki+2j2}(\{1,q^2\}^{ij+2})`$
$`=h_{mki+2j1}(\{1,q^2\}^{ij+2})+qh_{mki+2j2}(\{1,q^2\}^{ij+2}).`$
The result then follows from (28) and Theorem 3.1.
###### Corollary 3.3.
The polynomials $`P_{m,k}(q)`$ and $`Q_{m,k}(q)`$ have symmetric coefficients.
###### Proof.
A combinatorial way to see the symmetry of the coefficients of $`P_{m,k}(q)`$ is as follows: Modifying the weights such that vertical steps with an odd $`x`$-coordinate have weight $`q`$ and all the others weight $`1`$ does not change the entries of the determinant.
However, consider any given family of paths with weight $`q^w`$, when vertical steps with even $`x`$-coordinate have weight $`q`$. After the modification of the weights it will have weight $`q^{\mathrm{max}w}`$, where $`\mathrm{max}`$ is the total number of vertical steps in such a family of paths, which implies the claim.
For the polynomials $`Q_{m,k}`$, we use the following alternative weight: vertical steps with odd $`x`$-coordinate have weight $`q^2`$, vertical steps with starting point $`u_i`$ have weight $`1+q`$ and all others have weight $`1`$. ∎
When $`k=m1`$, there is only one lattice path from $`u_0=(0,0)`$ to $`v_0=(3,0)`$, which has weight $`1`$. This establishes the following result:
###### Corollary 3.4.
For $`m2`$, we have $`P_{m,m1}(q)=P_{m,m2}(q)`$ and $`Q_{m,m1}(q)=Q_{m,m2}(q)`$.
For the combinatorial interpretation of the $`q`$-Salié numbers, we need an auxiliary lemma:
###### Lemma 3.5.
Let $`(A_{ij})_{1i,jn}`$ and $`(B_{ij})_{1i,jn}`$ be two matrices. Then
$$\underset{1i,jn}{det}(A_{ij}+B_{ij})=\underset{I\{1,\mathrm{},n\}}{}\underset{1i,jn}{det}(D_{ij}^{(I)}),$$
where
$$D_{ij}^{(I)}=\{\begin{array}{cc}A_{ij},\hfill & \text{if }jI\text{,}\hfill \\ B_{ij},\hfill & \text{otherwise.}\hfill \end{array}$$
###### Theorem 3.6.
Let $`𝐮=(u_0,\mathrm{},u_{k1})`$ and $`𝐯=(v_0,\mathrm{},v_{k1})`$, where $`u_i=(2i,2i)`$ and $`v_i=(2i+2,mk1i)`$ for $`0ik1`$.
* The polynomial $`G_{m,k}(q)`$ is the sum of the weights of $`k`$-non-intersecting lattice paths $`𝐋`$ from $`𝐮`$ to $`𝐯`$ with the weight of $`𝐋`$ being
$$\underset{I\{0,1,\mathrm{},k1\}}{}w_I(𝐋),$$
where $`w_I`$ is defined as follows: for each $`iI`$, vertical steps with $`x`$-coordinate $`2i1`$ have weight $`q`$, and for any integer $`iI`$, vertical steps with $`x`$-coordinate $`2i`$ have weight $`q`$. All other steps have weight $`1`$.
* The polynomial $`H_{m,k}(q)`$ is the sum of the weights of $`k`$-non-intersecting lattice paths $`𝐋`$ from $`𝐮`$ to $`𝐯`$, with the weight of $`𝐋`$ being
$$\underset{I\{0,1,\mathrm{},k1\}}{}\overline{w}_I(𝐋),$$
where $`\overline{w}_I`$ is the same as $`w_I`$ – replacing $`q`$ with $`q^2`$ – with the exception of vertical steps starting from one of the points $`u_i`$, which have an additional weight of $`q`$. More precisely, if the weight of such a step would be $`1`$, it has weight $`1+q`$, it its weight would be $`q^2`$, it has weight $`q^2+q`$.
###### Proof.
(i) We apply Lemma 3.5 to $`det_{0i,jk1}(g_{mk+i+1,mk+j}(q))`$, where
$$\begin{array}{cc}\hfill g_{mk+i+1,mk+j}(q)& =h_{mki+2j1}(\{1\}^{ij+2},\{q\}^{ij+1})\hfill \\ & +h_{mki+2j1}(\{1\}^{ij+1},\{q\}^{ij+2}).\hfill \end{array}$$
Suppose that $`jI`$ and $`0ik1`$. Then we have to show that $`h_{mki+2j1}(\{1\}^{ij+2},\{q\}^{ij+1})`$ is the weighted sum of lattice paths from $`u_j`$ to $`v_i`$, where the vertical steps have the weight given in the claim. To this end, note that $`h_{mki+2j1}(\{1\}^{ij+2},\{q\}^{ij+1})`$ counts lattice paths from $`u_j`$ to $`v_i`$, when steps on $`ij+1`$ given vertical lines have weight $`q`$, those steps on the remaining $`ij+2`$ vertical lines have weight $`1`$.
By the construction in the claim, steps on exactly one of the vertical lines with $`x`$-coordinates $`2r1`$ and $`2r`$ have weight $`q`$. Since $`jI`$, steps on the vertical line with $`x`$-coordinate $`2j`$, i.e., with the $`x`$-coordinate of $`u_j`$, have weight $`1`$.
Similarly, if $`jI`$ we can verify that there are exactly $`ij+2`$ vertical lines between $`u_j`$ and $`v_i`$ with steps thereon having weight $`q`$.
(ii) In the same way, we can show that for $`jI`$ and $`0ik1`$.
$$h_{mki+2j1}(\{1\}^{ij+2},\{q^2\}^{ij+1})+qh_{mki+2j2}(\{1\}^{ij+2},\{q^2\}^{ij+1})$$
is the sum of weights of lattice paths from $`u_j`$ to $`v_i`$, where the vertical steps have the weight given in the claim. Meanwhile, for $`jI`$ and $`0ik1`$,
$$h_{mki+2j1}(\{1\}^{ij+1},\{q^2\}^{ij+2})+qh_{mki+2j2}(\{1\}^{ij+1},\{q^2\}^{ij+2})$$
is the sum of weights of lattice paths from $`u_j`$ to $`v_i`$. ∎
As an illustration of the underlying configurations in Theorem 3.6, we give an example in Figure 1 for $`m=7`$ and $`k=4`$.
###### Corollary 3.7.
The polynomials $`G_{m,k}(q)`$ and $`H_{m,k}(q)`$ have symmetric coefficients.
###### Proof.
A combinatorial way to see the symmetry of the coefficients of $`G_{m,k}(q)`$ is as follows: Modifying $`w_I`$ such that for each $`iI`$, vertical steps with $`x`$-coordinate $`2i`$ have weight $`q`$, and for any integer $`iI`$, vertical steps with $`x`$-coordinate $`2i1`$ have weight $`1`$ does not change the entries of the determinant.
However, consider any given family of paths with weight $`q^w`$ with weight by Theorem 3.6(i). After the modification of the weights it will have weight $`q^{\mathrm{max}w}`$, where $`\mathrm{max}`$ is the total number of vertical steps in such a family of paths, which implies the claim.
We omit the proof of the symmetry of the coefficients of $`H_{m,k}(q)`$. ∎
###### Corollary 3.8.
Let $`𝐮=(u_0,\mathrm{},u_{k1})`$ and $`𝐯=(v_0,\mathrm{},v_{k1})`$, where $`u_i=(2i,2i)`$ and $`v_i=(2i+2,mk1i)`$ for $`0ik1`$.
* The polynomial $`G_{m,k}(q)`$ is the sum of the weights of $`k`$-non-intersecting lattice paths $`𝐋`$ from $`𝐮`$ to $`𝐯`$ with the weight of $`𝐋`$ being
$$q^{\sigma _{2k}(𝐋)}\underset{i=0}{\overset{k1}{}}\left(q^{\sigma _{2i1}(𝐋)}+q^{\sigma _{2i}(𝐋)}\right),$$
where $`\sigma _j`$ denotes the number of vertical steps with $`x`$-coordinate $`j`$.
* The polynomial $`H_{m,k}(q)`$ is the sum of the weights of $`k`$-non-intersecting lattice paths $`𝐋`$ from $`𝐮`$ to $`𝐯`$ with the weight of $`𝐋`$ being
$$(1+q)^{f(𝐋)}q^{2\sigma _{2k}(𝐋)}\underset{i=0}{\overset{k1}{}}\left(q^{2\sigma _{2i1}(𝐋)}+q^{2\sigma _{2i}(𝐋)f_i(𝐋)}\right),$$
where $`\sigma _j`$ is as in (i) and $`f`$ (resp. $`f_i`$) denotes the number of vertical steps starting from $`𝐮`$ (resp. $`u_i`$).
###### Proof.
(i) By the definition of $`w_I`$, for $`0ik1`$, if $`iI`$, then vertical steps on the line with $`x`$-coordinates $`2i1`$ have weight $`q`$ and vertical steps on the line with $`x`$-coordinates $`2i`$ have weight $`1`$; and if $`iI`$, the case is just contrary. Note that steps on the vertical line with $`x`$-coordinates $`2k`$ always have weight $`q`$ and steps on the vertical line with $`x`$-coordinates $`2k1`$ always have weight $`1`$. This implies that
$$\underset{I\{0,1,\mathrm{},k1\}}{}w_I(𝐋)=q^{\sigma _{2k}(𝐋)}\underset{i=0}{\overset{k1}{}}\left(q^{\sigma _{2i1}(𝐋)}+q^{\sigma _{2i}(𝐋)}\right).$$
(ii) Notice that for $`0ik1`$, we have $`f_i(𝐋)=1`$ if $`𝐋`$ contains a vertical step starting from the point $`u_i`$, and $`f_i(𝐋)=0`$ otherwise. Similarly, we have
$`{\displaystyle \underset{I\{0,1,\mathrm{},k1\}}{}}\overline{w}_I(𝐋)`$ $`=q^{2\sigma _{2k}(𝐋)}{\displaystyle \underset{i=0}{\overset{k1}{}}}\left(q^{2\sigma _{2i1}(𝐋)}(1+q)^{f_i(𝐋)}+q^{2\sigma _{2i}(𝐋)2f_i(𝐋)}(q^2+q)^{f_i(𝐋)}\right),`$
$`=(1+q)^{f(𝐋)}q^{2\sigma _{2k}(𝐋)}{\displaystyle \underset{i=0}{\overset{k1}{}}}\left(q^{2\sigma _{2i1}(𝐋)}+q^{2\sigma _{2i}(𝐋)f_i(𝐋)}\right).`$
This completes the proof. ∎
The computation of $`G_{4,2}(q)`$ is illustrated in Figure 2, while the value of $`H_{4,2}(q)`$ as given in Table 4 is computed in Table 5.
###### Remark.
Since
$$\underset{1i,jn}{det}(A_{ij}+B_{ij})=\underset{I\{1,\mathrm{},n\}}{}\underset{1i,jn}{det}(C_{ij}^{(I)}),$$
where
$$C_{ij}^{(I)}=\{\begin{array}{cc}A_{ij},\hfill & \text{if }iI\text{,}\hfill \\ B_{ij},\hfill & \text{otherwise,}\hfill \end{array}$$
we may also define $`w_I`$ in Theorem 3.6(i) as follows: for each $`iI`$, vertical steps with $`x`$-coordinate $`2i+3`$ have weight $`q`$, and for any integer $`iI`$, vertical steps with $`x`$-coordinate $`2i+2`$ have weight $`q`$. All other steps have weight $`1`$. In this case, for each $`iI`$ and $`0jk1`$, we can show that $`h_{mki+2j1}(\{1\}^{ij+2},\{q\}^{ij+1})`$ is the weighted sum of lattice paths from $`u_j`$ to $`v_i`$. Moreover,
$$\underset{I\{0,1,\mathrm{},k1\}}{}w_I(𝐋)=q^{\sigma _0(𝐋)}\underset{i=1}{\overset{k}{}}\left(q^{\sigma _{2i}(𝐋)}+q^{\sigma _{2i+1}(𝐋)}\right).$$
Similarly, we may define $`\overline{w}_I`$ in Theorem 3.6(ii) as follows: for each $`iI`$, a vertical step toward the point $`v_i`$ has weight $`q+1`$, vertical steps with $`x`$-coordinate $`2i+3`$ have weight $`q^2`$. For any integer $`iI`$, a vertical step toward the point $`v_i`$ has weight $`q^2+q`$, and vertical steps with $`x`$-coordinate $`2i+2`$ not toward $`v_i`$ have weight $`q^2`$. All other steps have weight $`1`$. In this case, we have
$$\underset{I\{0,1,\mathrm{},k1\}}{}\overline{w}_I(𝐋)=(1+q)^{\overline{f}(𝐋)}q^{2\sigma _0(𝐋)}\underset{i=1}{\overset{k}{}}\left(q^{2\sigma _{2i}(𝐋)\overline{f}_i(𝐋)}+q^{2\sigma _{2i+1}(𝐋)}\right),$$
where $`\overline{f}`$ (resp. $`\overline{f}_i`$) denotes the number of vertical steps ending in $`𝐯`$ (resp. $`v_i`$).
When $`k=m1`$, there is only one lattice path from $`u_0=(0,0)`$ to $`v_0=(2,0)`$, which has weight $`1`$. This establishes the following result:
###### Corollary 3.9.
$`G_{m,m1}(q)=2G_{m,m2}(q)`$ and $`H_{m,m1}(q)=2H_{m,m2}(q)`$.
## 4. Open problems
We would like to point out three directions of possible further research: It appears that the polynomials $`P_{m,k}`$ and $`G_{m,k}`$ are log-concave, however, we did not pursue this question further. Note that the polynomials $`Q_{m,k}`$ and $`H_{m,k}`$ are not even unimodal.
Victor Guo and Jiang Zeng gave in even finer $`q`$-analogues of the polynomials considered here, replacing (3) and (4) by
$`S_{m,n,r}(q)`$ $`={\displaystyle \underset{k=1}{\overset{n}{}}}{\displaystyle \frac{[2rk]}{[2r]}}[k]^{m1}q^{\frac{m+2r1}{2}(nk)},`$
$`T_{m,n,r}(q)`$ $`={\displaystyle \underset{k=1}{\overset{n}{}}}(1)^{nk}{\displaystyle \frac{[(2r1)k]}{[2r1]}}[k]^{m1}q^{\frac{m}{2}(nk)},`$
where $`r1`$.
Although the coefficients of the corresponding polynomials $`P_{m,k,r},Q_{m,k,r},G_{m,k,r}`$ and $`H_{m,k,r}`$ are not positive anymore, one might hope for a refinement of Theorem 2.3.
Finally, we should point out that Ira Gessel and Xavier Viennot also presented nice generating functions for their coefficients $`f_{m,k}`$ and $`s_{m,k}`$, namely
$`{\displaystyle \underset{m,k}{}}s_{m,k}t^k{\displaystyle \frac{x^{2n}}{(2n)!}}`$ $`={\displaystyle \frac{\mathrm{cosh}\sqrt{1+4t}\frac{x}{2}}{\mathrm{cosh}\frac{x}{2}}},`$
$`{\displaystyle \underset{m,k}{}}f_{m,k}t^k{\displaystyle \frac{x^{2n+1}}{(2n+1)!}}`$ $`={\displaystyle \frac{\mathrm{cosh}\sqrt{1+4t}\frac{x}{2}\mathrm{cosh}\frac{x}{2}}{t\mathrm{sinh}\frac{x}{2}}}.`$
It would be interesting to find the corresponding refinements.
## 5. Epilogue
One may wonder how these results were discovered. The truth is, that at first “only” formula (5) was known. Using this formula, Table 1 was computed. Then, in analogy to , the matrix
$$\left((1)^{km}\frac{[m]!}{[k+1]!}P_{k,km}(q)\right)_{0k,mn}$$
was inverted and, since we were looking for a lattice path interpretation, the entry in row $`i`$ and column $`j`$ of the inverse matrix had to be the weighted number of lattice paths from $`u_j`$ to $`v_i`$. This given, it was easy to find the correct weights. Finally, we read the proof given in backwards, its first line corresponding to our Lemma 2.2.
Acknowledgment. The third author was supported by EC’s IHRP Programme, within Research Training Network “Algebraic Combinatorics in Europe,” grant HPRN-CT-2001-00272.
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# Photon-propagation model with random background field: Length scales and Cherenkov limits
## I Introduction
In a previous article KlinkhamerRupp , we have proposed a simple photon-propagation model to describe the potential effects of a static spacetime foam composed of identical, randomly-distributed defects (e.g., microscopic wormholes) embedded in Minkowski spacetime. For this particular model, a modified photon dispersion law was derived in the long-wavelength limit,
$$\omega ^2(1A^2\gamma _1)c^2k^2A^2l_\gamma ^2c^2k^4,$$
(1)
where $`k|\stackrel{}{k}|`$ is the photon wave number and $`\omega `$ the frequency, $`A`$ the amplitude of the frozen (time-independent) random background field $`g_1(\stackrel{}{x})`$, $`l_\gamma `$ a characteristic length scale of $`g_1(\stackrel{}{x})`$, $`\gamma _1`$ a nonnegative dimensionless coefficient, and $`c`$ a fundamental constant tracing back to the Minkowski line element (see also below). An upper bound $`l_\gamma <1.6\times 10^{22}\mathrm{cm}`$, for $`A=\alpha 1/137`$, was then obtained from observations of a particular TeV flare in an active galactic nucleus.
In this Brief Report, we use recent results on ultra-high-energy cosmic rays ColemanGlashow ; GagnonMoore to improve our previous bound on $`l_\gamma `$. In addition, we give a careful discussion of the possible relation between the photonic length scale $`l_\gamma `$ and the characteristic length scales of the microscopic spacetime structure.
In the following, we will use standard natural units with $`\mathrm{}`$ $`=`$ $`c`$ $`=`$ $`1`$, except when stated otherwise. For the physical situation discussed in the next section, the operational definition of the velocity $`c`$ is the maximum attainable velocity of the proton. (Further discussions on Lorentz noninvariance can be found in, e.g., Refs. ColemanGlashow ; GagnonMoore ; ColladayKostelecky ; Bertolami-etal and references therein.)
## II Photon propagation
Assuming a modified photon dispersion law with a negative dimensionful coefficient $`K_{1\mathrm{neg}}`$,
$$E_\gamma k+K_{1\mathrm{neg}}k^3,$$
(2)
and an unchanged (ultrarelativistic) proton dispersion law,
$$E_pk,$$
(3)
the Cherenkov-like proton process $`pp+\gamma `$ becomes kinematically allowed ColemanGlashow . From observations of ultra-high-energy cosmic rays, Gagnon and Moore GagnonMoore obtain the following bound:
$$0K_{1\mathrm{neg}}<(4\times 10^{22}\mathrm{GeV})^2.$$
(4)
There is also a bound on the difference between the maximum attainable velocities of particles with spin $`1`$ and spin $`1/2`$. For a modified photon dispersion law
$$E_\gamma =c(1+ϵ)k$$
(5)
and an unmodified fermion dispersion law (3), the authors of Ref. GagnonMoore obtain the bound
$$|ϵ|<1.6\times 10^{23}.$$
(6)
We now turn to a simple photon-propagation model KlinkhamerRupp with a fixed random background field $`g_1(x)`$ and an action given by
$`S_{\mathrm{photon}}=`$ $`\frac{1}{4}_^4\mathrm{d}^4x\left(F_{\mu \nu }(x)F_{\kappa \lambda }(x)\eta ^{\kappa \mu }\eta ^{\lambda \nu }+g_1(x)F_{\kappa \lambda }(x)\stackrel{~}{F}^{\kappa \lambda }(x)\right),`$ (7)
in terms of the standard Maxwell field strength tensor $`F_{\mu \nu }_\mu A_\nu _\nu A_\mu `$, the dual tensor $`\stackrel{~}{F}^{\kappa \lambda }\frac{1}{2}ϵ^{\kappa \lambda \mu \nu }F_{\mu \nu }`$, and the inverse Minkowski metric $`\eta ^{\mu \nu }`$. The random background field $`g_1(x)`$ is assumed to be time-independent,
$$g_1(x^0,x^1,x^2,x^3)=g_1(x^1,x^2,x^3)g_1(\stackrel{}{x}),$$
(8)
and to fluctuate around a value zero with amplitude $`A`$; see Sec. IV of Ref. KlinkhamerRupp for further properties. The random background field $`g_1(\stackrel{}{x})`$ in Eq. (7) can be seen to act as a variable coupling constant, with spacetime taken to be perfectly smooth (manifold $`\mathrm{M}=^4`$ and metric $`g_{\mu \nu }(x)=\eta _{\mu \nu })`$.
For the photon-propagation model (7), we have calculated in Sec. V of Ref. KlinkhamerRupp the dispersion law (1), with $`l_\gamma `$ and $`\gamma _1`$ determined in terms of the autocorrelation function of $`g_1(\stackrel{}{x})`$,
$$l_\gamma =l_\gamma \left[g_1\right],\gamma _1=\gamma _1\left[g_1\right].$$
(9)
The dispersion law (1) gives then the following photon energy:
$$E_\gamma k\left(1\frac{1}{2}A^2\gamma _1\frac{1}{2}A^2l_\gamma ^2k^2\right),$$
(10)
for parametrically small amplitude $`A`$ or for $`\gamma _1`$ and $`l_\gamma ^2k^2`$ much less than unity.
From the experimental bounds (4), (6) and the relation (10), we obtain
$`l_\gamma `$ $`<\left(2\times 10^{20}\mathrm{GeV}\right)^1\left(\alpha A^1\right)\left(1.0\times 10^{34}\mathrm{cm}\right)\left(1/137A^1\right),`$ (11a)
$`\gamma _1`$ $`<\left(6\times 10^{19}\right)\left(\alpha /A\right)^2,`$ (11b)
where $`\alpha 1/137`$ is the fine-structure constant (the possible relation $`A\alpha `$ will be discussed in the next section). Note that the bound (11a) is $`12`$ orders of magnitude better than the one given in Sec. VI of Ref. KlinkhamerRupp , where $`l_\gamma `$ was called $`l_{\mathrm{foam}}`$. This bound on $`l_\gamma `$ for $`A\alpha `$ is, in fact, of the order of the Planck length,
$$l_{\mathrm{Planck}}\sqrt{G\mathrm{}/c^3}1.6\times 10^{33}\mathrm{cm},$$
(12)
which may determine the fine-scale structure of spacetime itself Wheeler .
## III Spacetime structure
In order to connect the photon parameters $`\gamma _1\left[g_1\right]`$ and $`l_\gamma \left[g_1\right]`$ derived from the effective action (7) to the microscopic structure of spacetime, we introduce the following definitions:
$`l_\gamma `$ $`l_{\mathrm{wormhole}}\left(l_{\mathrm{wormhole}}/l_{\mathrm{separation}}\right)^{3/2},`$ (13a)
$`\gamma _1`$ $`\left(l_{\mathrm{wormhole}}/l_{\mathrm{separation}}\right)^3.`$ (13b)
These definitions are motivated by a very simple spacetime model KlinkhamerRupp consisting of static, randomly-distributed wormholes Wheeler embedded in Minkowski spacetime. This toy model has, by definition, a preferred frame of reference. The length $`l_{\mathrm{wormhole}}`$ would then correspond to an appropriate characteristic dimension of an *individual* wormhole (e.g., the average width of the two mouths or the long distance between the centers of the mouths, where both lengths are measured in the Minkowski part of spacetime and the short distance through the wormhole throat is assumed to be zero). The length $`l_{\mathrm{separation}}`$ would correspond to the average separation between *different* wormholes (the wormhole density is $`n_{\mathrm{wormholes}}=l_{\mathrm{separation}}^3`$).
The anomaly calculation reported in the Appendix of Ref. KlinkhamerRupp , specialized to the case $`l_h=\delta `$ and with notations $`(l_{\mathrm{foam}},d,a)`$ for $`(l_\gamma ,l_{\mathrm{wormhole}},l_{\mathrm{separation}})`$ here, gave $`A=\alpha `$ in the effective action (7) and extra factors $`0.18`$ and $`0.15`$ on the right-hand sides of Eqs. (13a) and (13b), respectively. This calculation was, however, based on several simplifying assumptions and is, therefore, not absolutely rigorous. The two most important results would be that there are no extremely small or large factors on the right-hand sides of Eqs.(13ab) and that the effective amplitude $`A`$ is of order $`\alpha `$. The physical interpretation of the quantities $`l_{\mathrm{wormhole}}`$ and $`l_{\mathrm{separation}}`$, defined mathematically by Eqs. (13ab), would be that they emerge directly from the underlying spacetime structure. Indeed, a successful calculation would relate the “randomness” of the couplings $`g_1(\stackrel{}{x})`$ in the effective action (7) to the (as of yet, unknown) microscopic structure of spacetime. (An entirely different origin for the variable couplings $`g_1(x)`$ of a $`F_{\kappa \lambda }\stackrel{~}{F}^{\kappa \lambda }`$ term in the effective action is, of course, not excluded; see, e.g., Ref. Bertolami-etal and references therein.)
A concrete example of this particular spacetime model with permanent wormholes would then have
$$l_{\mathrm{wormhole}}10\times l_{\mathrm{Planck}},l_{\mathrm{separation}}10^8\times l_{\mathrm{Planck}},$$
(14)
in order to be consistent with the bounds (11a) and (11b) for $`A=\alpha `$. More generally, Fig. 1 shows which combinations of values of $`l_{\mathrm{wormhole}}`$ and $`l_{\mathrm{separation}}`$ are allowed or excluded, assuming $`A=\alpha `$. For $`l_{\mathrm{wormhole}}1.3\times 10^{25}\mathrm{cm}`$ (or $`l_{\mathrm{separation}}1.5\times 10^{19}\mathrm{cm}`$) the bound (11b) is seen to be the stronger one and for the other case the bound (11a). Without further input, we cannot say anything about $`l_{\mathrm{wormhole}}`$ and $`l_{\mathrm{separation}}`$ individually.
## IV Discussion
Using experimental bounds GagnonMoore on possible Lorentz-violating modifications of the photon dispersion law from the absence of Cherenkov-like processes for high-energy cosmic rays, we have obtained bounds on the length scales of a photon-propagation model (7) with time-independent random background field, which could result from a static, multiply connected spacetime foam KlinkhamerRupp . Even though the effective length scale $`l_\gamma `$ which enters the photon dispersion law is constrained to be below the Planck length $`l_{\mathrm{Planck}}`$ for $`A=\alpha `$, these bounds do not rigorously exclude a foamlike structure of spacetime with length scales $`l_{\mathrm{wormhole}}`$ and $`l_{\mathrm{separation}}`$ at or even above the Planck length (see Fig. 1).
On the other hand, it would perhaps not be unreasonable to expect Wheeler some remnant “quantum-gravity” effect with *both* length scales $`l_{\mathrm{wormhole}}`$ and $`l_{\mathrm{separation}}`$ of the order of the Planck length (12), even for a time-independent model with corresponding preferred frame of reference. But the static wormhole gas with $`l_{\mathrm{wormhole}}l_{\mathrm{separation}}l_{\mathrm{Planck}}10^{33}\mathrm{cm}`$ and $`A\alpha `$ is ruled out by the bounds (11ab) in terms of (13ab); cf. Fig. 1. \[The crucial assumption here is that the (static) spacetime foam gives rise to an effective theory (7) with $`g_1(\stackrel{}{x})`$ amplitude $`A`$ of order $`\alpha `$. If, for some reason, $`A`$ would be very much smaller than $`\alpha `$, the bounds (11ab) become essentially inoperative. As mentioned in the previous section, the preliminary calculations of Ref. KlinkhamerRupp do suggest $`A\alpha `$, but this remains to be confirmed.\]
The tentative conclusion is, therefore, that a preferred-frame graininess of space with a single length scale $`l_{\mathrm{Planck}}`$ may be hard to reconcile with the current experimental bounds from cosmic-ray physics. Without fine-tuning, such a graininess of space can also be expected Collins-etal to show up in “low-energy” physics (i.e., $`\sqrt{s}E_{\mathrm{Planck}}\sqrt{\mathrm{}c^5/G}1.2\times 10^{19}\mathrm{GeV}`$) with powers of the coupling constants as the only suppression factor, an example being the linear term of Eq. (10) with $`A^2\gamma _1\alpha ^2`$. One possible solution would have gravity as an emergent phenomenon and the Lorentz-violation scale moved to trans-Planckian energies KlinkhamerVolovik . But, this is only one out of many suggestions and the puzzle of the apparent smoothness of space remains unsolved.
## ACKNOWLEDGMENTS
FRK thanks G.E. Volovik for interesting discussions and, also, S. Liberati, V.A. Rubakov, and the other participants of the Conference on Fundamental Symmetries and Fundamental Constants (ICTP, Trieste, September 2004). CR thanks A. Mazumdar for useful comments after a seminar (NBI, Copenhagen, April 2005).
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# Thermalization vs. Isotropization & Azimuthal Fluctuations11footnote 1presented at the workshop Correlations and Fluctuations in Relativistic Nuclear Collisions, MIT, April 21-23, 2005
## 1 Introduction
A matter created in relativistic heavy-ion collisions manifests a strongly collective hydrodynamic behaviour , particularly evident in studies of the so-called elliptic flow . Hydrodynamic description requires, strictly speaking, a local thermal equilibrium and experimental data on the particle spectra and the elliptic flow suggest, when analysed within the hydrodynamic model, that an equilibration time of the parton<sup>2</sup><sup>2</sup>2The term ‘parton’ is used to denote a fermionic (quark) or bosonic (gluon) excitation of the quark-gluon plasma. system produced at the collision early stage is as short as 0.6 $`\mathrm{fm}/c`$ . Such a fast equilibration can be explained assuming that the quark-gluon plasma is strongly coupled . However, high-energy density in the collision early stage, when the elliptic flow is generated , allows one to believe that the plasma is then weakly coupled due to the asymptotic freedom. Calculations, which assume that the parton-parton collisions are responsible for the equilibration of the weakly interacting plasma, provide an equilibration time of at least 2.6 $`\mathrm{fm}/c`$ . To thermalize the system one needs either a few hard collisions of the momentum transfer of order of the characteristic parton momentum<sup>3</sup><sup>3</sup>3Although I consider anisotropic systems, the characteristic momentum in all directions is assumed to be of the same order., which we denote here as $`T`$ (as the temperature of equilibrium system), or many collisions of smaller transfer. As discussed in e.g. , the inverse time scale of the collisional equilibration is of order $`g^4\mathrm{ln}(1/g)T`$ where $`g`$ is the QCD coupling constant. However, it has been argued that the equilibration is speeded up by instabilities generated in an anisotropic quark-gluon plasma , as growth of the unstable modes is associated with the system’s isotropization. The characteristic inverse time of instability development is roughly of order $`gT`$ for a sufficiently anisotropic momentum distribution . Thus, the instabilities are much ‘faster’ than the hard collisions in the weak coupling regime. Very recent classical simulation indeed shows effectiveness of the instabilities driven isotropization.
The isotropization should be clearly distinguished from the equilibration process . The instabilities driven isotropization is a mean-field reversible phenomenon which is not accompanied with entropy production. Therefore, the collisions, which are responsible for the dissipation, are needed to reach the equilibrium state of maximal entropy. The instabilities contribute to the equilibration indirectly, reducing relative parton momenta and increasing the collision rate.
It has been recently observed that the hydrodynamic collective behaviour does not actually require local thermodynamic equilibrium but a merely isotropic momentum distribution of liquid components . Thus, there is a question whether a quark-gluon plasma, which is equilibrated nearly immediately after its production as advocated in , can be distinguished from the parton system which slowly evolves towards equilibrium being isotropized fast. I argue here that measurements of azimuthal fluctuations, which are generated at the early stage of heavy-ion collisions, can help to distinguish the two scenarios.
In the first part of my talk I review the instabilities driven isotropization. I discuss how the unstable modes are initiated and what is a mechanism responsible for their growth. Dispersion relations of the unstable modes are considered, and it is explained why the development of instabilities is associated with the system’s isotropization. In the second part of my talk I discuss the azimuthal fluctuations, arguing that the fluctuations generated in the non-equilibrium isotropic system are much larger than those in the fully equilibrated plasma. Two possible measurements are proposed.
## 2 Instabilities driven isotropization
Temporal evolution of the electron-ion plasma is plagued by a large variety of instabilities. Those caused by coordinate space inhomogeneities, in particular by the system’s boundaries, are usually called the hydrodynamic instabilities while those due to non-equilibrium momentum distribution of plasma particles the kinetic instabilities. Hardly anything is known about hydrodynamic instabilities of the quark-gluon plasma, and I will not speculate about their possible role in the system’s dynamics. The kinetic instabilities are initiated either by the charge or current fluctuations. In the first case, the electric field ($`𝐄`$) is longitudinal ($`𝐄𝐤`$, where $`𝐤`$ is the wave vector), while in the second case the field is transverse ($`𝐄𝐤`$). For this reason, the kinetic instabilities caused by the charge fluctuations are usually called longitudinal while those caused the current fluctuations transverse. Since the electric field plays a crucial role in the longitudinal mode generation, the longitudinal instabilities are also called electric while the transverse ones magnetic.
In the non-relativistic plasma the electric instabilities are usually much more important than the magnetic ones as the magnetic effects are suppressed by the factor $`v^2/c^2`$ where $`v`$ is the particle’s velocity. In the relativistic plasma both types of similar strength. The electric instabilities occur when the momentum distribution has more than one maximum while a sufficient condition for the magnetic instabilities is, as discussed below, anisotropy of the momentum distribution. For this reason, the magnetic unstable mode, which is also called Weibel or filamentation instability , was argued long ago to be relevant for equilibration of the quark-gluon plasma produced in relativistic heavy-ion collisions . In the remaing part of the section, I am going to explain in detail why the filamentation is relevant and how it speeds up the process of plasma thermalization.
### 2.1 Seeds of the filamentation
Let me first discuss how the unstable transverse modes are initiated. For this purpose I consider a parton system which is homogeneous but the parton momentum distribution is, in general, not of the equilibrium form, it is not isotropic. The system is on average locally colourless but colour fluctuations are possible. Therefore, $`j_a^\mu (x)=0`$ where $`j_a^\mu (x)`$ is a local colour four-current in the adjoint representation of $`\mathrm{SU}(3)`$ gauge group with $`\mu =0,1,2,3`$ and $`a=1,2,3,\mathrm{},8`$ being the Lorentz and colour index, respectively; $`x=(t,𝐱)`$ denotes a four-position in the coordinate space.
As discussed in detail in , the current correlator for a classical system of non-interacting quarks and gluons is
$`M_{ab}^{\mu \nu }(t,𝐱)`$ $`\stackrel{\mathrm{def}}{=}`$ $`j_a^\mu (t_1,𝐱_1)j_b^\nu (t_2,𝐱_2)={\displaystyle \frac{1}{8}}g^2\delta ^{ab}{\displaystyle \frac{d^3p}{(2\pi )^3}\frac{p^\mu p^\nu }{E_p^2}f(𝐩)\delta ^{(3)}(𝐱𝐯t)},`$ (1)
where $`(t,𝐱)(t_2t_1,𝐱_2𝐱_1)`$ and the effective parton distribution function $`f(𝐩)`$ equals $`n(𝐩)+\overline{n}(𝐩)+6n_g(𝐩)`$ with $`n(𝐩)`$, $`\overline{n}(𝐩)`$ and $`n_g(𝐩)`$ giving the average colourless distribution function of quarks $`Q^{ij}(x,𝐩)=\delta ^{ij}n(𝐩)`$, antiquarks $`\overline{Q}^{ij}(x,𝐩)=\delta ^{ij}\overline{n}(𝐩)`$, and gluons $`G^{ab}(x,𝐩)=\delta ^{ab}n_g(𝐩)`$. We note that the distribution function of (anti-)quarks belongs to the fundamental representation of the $`\mathrm{SU}(3)`$ gauge while that of gluons to the adjoint representation. Therefore, $`i,j=1,2,3`$ and $`a,b=1,2,\mathrm{},8`$.
Due to the average space-time homogeneity, the correlation tensor (1) depends only on the difference $`(t_2t_1,𝐱_2𝐱_1)`$. The space-time points $`(t_1,𝐱_1)`$ and $`(t_2,𝐱_2)`$ are correlated in the system of non-interacting particles if a particle travels from $`(t_1,𝐱_1)`$ to $`(t_2,𝐱_2)`$. For this reason the delta $`\delta ^{(3)}(𝐱𝐯t)`$ is present in the formula (1). The momentum integral of the distribution function simply represents the summation over particles. The fluctuation spectrum is found as a Fourier transform of the tensor (1) i.e.
$$M_{ab}^{\mu \nu }(\omega ,𝐤)=\frac{1}{8}g^2\delta ^{ab}\frac{d^3p}{(2\pi )^3}\frac{p^\mu p^\nu }{E_p^2}f(𝐩)\mathrm{\hspace{0.33em}2}\pi \delta (\omega \mathrm{𝐤𝐯}).$$
(2)
To compute the fluctuation spectrum, the parton momentum distribution has to be specified. Such calculations with two forms of the anisotropic momentum distribution are presented in . Here I only qualitatively discuss Eqs. (1,2). I assume that the momentum distribution is elongated in, say, the $`z`$ direction. Then, Eqs. (1,2) clearly show that the corelator $`M^{zz}`$ is larger than $`M^{xx}`$ or $`M^{yy}`$. It also clear that $`M^{zz}`$ is the largest when the wave vector $`𝐤`$ is along the direction of the momentum deficit. Then, the delta function $`\delta (\omega \mathrm{𝐤𝐯})`$ does not much constrain the integral in Eq. (2). Since the momentum distribution is elongated in the $`z`$ direction, the current fluctuations are the largest when the wave vector $`𝐤`$ is the $`xy`$ plane. Thus, I conclude that some fluctuations in the anisotropic system are large, much larger than in the isotropic one and that anisotropic system has a natural tendency to split into the current filaments parallel to the direction of the momentum surplus. These currents are seeds of the filamentation instability.
### 2.2 Mechanism of filamentation
Let me now explain in terms of elementary physics why the fluctuating currents, which flow in the direction of the momentum surplus, can grow in time. To simplify the discussion, which follows , I consider an electromagnetic anisotropic system. The form of the fluctuating current is chosen to be
$$𝐣(x)=j\widehat{𝐞}_z\mathrm{cos}(k_xx),$$
(3)
where $`\widehat{𝐞}_z`$ is the unit vector in the $`z`$ direction. As seen in Eq. (3), there are current filaments of the thickness $`\pi /|k_x|`$ with the current flowing in the opposite directions in the neighbouring filaments.
The magnetic field generated by the current (3) is given as
$`𝐁(x)={\displaystyle \frac{j}{k_x}}\widehat{𝐞}_y\mathrm{sin}(k_xx),`$
and the Lorentz force acting on the partons, which fly along the $`z`$ direction, equals
$`𝐅(x)=q𝐯\times 𝐁(x)=qv_z{\displaystyle \frac{j}{k_x}}\widehat{𝐞}_x\mathrm{sin}(k_xx),`$
where $`q`$ is the electric charge. One observes, see Fig. 1, that the force distributes the partons in such a way that those, which positively contribute to the current in a given filament, are focused in the filament centre while those, which negatively contribute, are moved to the neighbouring one. Thus, the initial current is growing.
### 2.3 Dispersion equation
The Fourier transformed chromodynamic field $`A^\mu (k)`$ satisfies the equation of motion as
$$\left[k^2g^{\mu \nu }k^\mu k^\nu \mathrm{\Pi }^{\mu \nu }(k)\right]A_\nu (k)=0,$$
(4)
where $`\mathrm{\Pi }^{\mu \nu }(k)`$ is the polarization tensor or gluon self-energy which is discussed later on. A general plasmon dispersion equation is of the form
$$\mathrm{det}\left[k^2g^{\mu \nu }k^\mu k^\nu \mathrm{\Pi }^{\mu \nu }(k)\right]=0.$$
(5)
Equivalently, the dispersion relations are given by the positions of poles of the effective gluon propagator. Due to the transversality of $`\mathrm{\Pi }^{\mu \nu }(k)`$ ($`k_\mu \mathrm{\Pi }^{\mu \nu }(k)=k_\nu \mathrm{\Pi }^{\mu \nu }(k)=0`$) not all components of $`\mathrm{\Pi }^{\mu \nu }(k)`$ are independent from each other, and consequently the dispersion equation (5), which involves a determinant of $`4\times 4`$ matrix, can be simplified to the determinant of $`3\times 3`$ matrix. For this purpose I introduce the colour permittivity tensor $`ϵ^{lm}(k)`$ where the indices $`l,m,n=1,2,3`$ label three-vector and tensor components. Because of the relation
$$ϵ^{lm}(k)E^l(k)E^m(k)=\mathrm{\Pi }^{\mu \nu }(k)A_\mu (k)A_\nu (k),$$
where $`𝐄`$ is the chromoelectric vector, the permittivity can be expressed through the polarization tensor as
$$ϵ^{lm}(k)=\delta ^{lm}+\frac{1}{\omega ^2}\mathrm{\Pi }^{lm}(k).$$
Then, the dispersion equation gets the form
$$\mathrm{det}\left[𝐤^2\delta ^{lm}k^lk^m\omega ^2ϵ^{lm}(k)\right]=0.$$
(6)
The relationship between Eq. (5) and Eq. (6) is most easily seen in the Coulomb gauge when $`A^0=0`$ and $`𝐤𝐀(k)=0`$. Then, $`𝐄=i\omega 𝐀`$ and Eq. (4) is immediately transformed into an equation of motion of $`𝐄(k)`$ which further provides the dispersion equation (6).
The dynamical information is contained in the polarization tensor $`\mathrm{\Pi }^{\mu \nu }(k)`$ or, equivalently, in the permittivity tensor $`ϵ^{lm}(k)`$ which can be derived either within the transport theory or diagrammatically . The result is
$$ϵ^{nm}(\omega ,𝐤)=\delta ^{nm}+\frac{g^2}{2\omega }\frac{d^3p}{(2\pi )^3}\frac{v^n}{\omega \mathrm{𝐤𝐯}+i0^+}\frac{f(𝐩)}{p^l}\left[\left(1\frac{\mathrm{𝐤𝐯}}{\omega }\right)\delta ^{lm}+\frac{k^lv^m}{\omega }\right].$$
(7)
Since $`\mathrm{\Pi }^{\mu \nu }(k)`$ and $`ϵ^{lm}(k)`$ are unit matrices in the colour space, the colour indices are suppressed here.
Substituting the permittivity (7) into Eq. (6), one fully specifies the dispersion equation (6) which provides a spectrum of quasi-particle bosonic excitations. A solution $`\omega (𝐤)`$ of Eq. (6) is called stable when $`\mathrm{Im}\omega 0`$ and unstable when $`\mathrm{Im}\omega >0`$. In the first case the amplitude is constant or it exponentially decreases in time while in the second one there is an exponential growth of the amplitude. In practice it appears difficult to find solutions of Eq. (6) because of rather complicated structure of the tensor (7). However, the problem simplifies as we are interested in specific modes which are expected to be unstable. Namely, we look for solutions corresponding to the fluctuating current in the direction of the momentum surplus and the wave vector perpendicular to it.
As previously, the momentum distribution is assumed to be elongated in the $`z`$ direction, and consequently the fluctuating current also flows in this direction. The magnetic field has a non-vanishing component along the $`y`$ direction and the electric filed in the $`z`$ direction. Finally, the wave vector is parallel to the axis $`x`$, see Fig. 1. We also assume that the momentum distribution obeys the mirror symmetry $`f(𝐩)=f(𝐩)`$, and then the permittivity tensor has only non-vanishing diagonal components. Taking into account all these conditions, one simplifies the dispersion equation (6) to the form
$$H(\omega )k_x^2\omega ^2ϵ^{zz}(\omega ,k_x)=0,$$
(8)
where only one diagonal component of the dielectric tensor enters.
It appears that an existence of unstable solutions of Eq. (8) can be proved without solving it. The so-called Penrose criterion , which follows from analytic properties of the permittivity as a function of $`\omega `$, states that the dispersion equation $`H(\omega )=0`$ has unstable solutions if $`H(\omega =0)<0`$. The Penrose criterion was applied to the equation (8) in but a much more general discussion of the instability condition is presented in . Not entering into details, there exist unstable modes if the momentum distribution averaged (with a proper weight) over momentum length is anisotropic.
To solve the dispersion equation (8), the parton momentum distribution has to be specified. Several analytic (usually approximate) solution of the dispersion equation with various momentum distributions can be found in . An example of the numerical solution, which gives the unstable mode frequency in the full range of wave vectors is shown Fig. 2 taken from . The momentum distribution is of the form
$$f(𝐩)\frac{1}{(p_T^2+\sigma _{}^2)^3}\mathrm{e}^{\frac{p_z^2}{2\sigma _{}^2}},$$
where $`p_{}\sqrt{p_x^2+p_y^2}`$. The mode is pure imaginary and $`\gamma _k\mathrm{Im}\omega (k_{})`$. The value of the coupling is $`\alpha _sg^2/4\pi =0.3`$, $`\sigma _{}=0.3\mathrm{GeV}`$ and the effective parton density is chosen to be $`6\mathrm{fm}^3`$. As seen, there is a finite interval of wave vectors for which the unstable modes exist.
### 2.4 Growth of instabilities and abelianization
A time evolution of a classical many-parton system interacting via classical chromodynamic field has been studied in . Numerical simulations have been performed effectively in $`1+1`$ dimensions as the chromodynamic potentials depend on $`t`$ and $`x`$. The initial field amplitudes are assumed to obey Gaussian white noise and the initial parton momentum distribution is
$$f(𝐩)\delta (p_z)\mathrm{e}^{\frac{\sqrt{p_x^2+p_y^2}}{p_{\mathrm{hard}}}},$$
(9)
with $`p_{\mathrm{hard}}=10\mathrm{GeV}`$.
Fig. 3, which is taken from , shows results of the simulation corresponding to a lattice of physical size $`L=40\mathrm{fm}`$. As seen, the amount of energy of the fields, which is initially much smaller than the kinetic energy of all particles, grows exponentially and the magnetic contribution dominates. The simulation indirectly confirms existence of the unstable magnetic modes in the system.
Unstable modes cannot grow to infinity and the question arises what is a mechanism responsible for stopping the instability growth. One suspects that non-Abelian non-linearities can play an important role here. An elegant qualitative argument suggests that the non-linearities do not stabilize the unstable modes because the system spontaneously chooses an Abelian configuration in the course of the instability development. Let me explain the idea.
In the Coulomb gauge the effective potential of the unstable configuration has the form
$`V_{\mathrm{eff}}[𝐀^a]=\mu ^2𝐀^a𝐀^a+{\displaystyle \frac{1}{4}}g^2f^{abc}f^{ade}(𝐀^b𝐀^d)(𝐀^c𝐀^e),`$
which is shown in Fig. 4 taken from . The first term (with $`\mu ^2>0`$) is responsible for a very existence of the instability. The second term, which comes from the Yang-Mills lagrangian, is of pure non-Abelian nature. The term appears to be positive and thus it counteracts the instability growth. However, the non-Abelian term vanishes when the potential $`𝐀^a`$ is effectively Abelian, and consequently, such a configuration corresponds to the steepest decrease of the effective potential. Thus, the system spontaneously abelianizes in the course of instability growth.
The effect of abelianization has been indeed found in numerical simulations performed in the $`1+1`$ dimensions . As an example, I show in Fig. 5 the result of fully classical simulation . One observes in Fig. 5 taken from , where
$`\varphi _{\mathrm{rms}}\sqrt{{\displaystyle _0^L}{\displaystyle \frac{dx}{L}}(A_y^aA_y^a+A_z^aA_z^a)},\overline{C}{\displaystyle _0^L}{\displaystyle \frac{dx}{L}}{\displaystyle \frac{\sqrt{\mathrm{Tr}[(i[A_y,A_z])^2]}}{\mathrm{Tr}[A_y^2+A_z^2]}},`$
that the field commutator measured by $`\overline{C}`$ decreases in time, in spite of the field growth quantified by $`\varphi _{\mathrm{rms}}`$.
Very recent simulations performed in the $`1+3`$ dimensions , which utilise a complete Hard Loop action for anisotropic systems , show that the growth of unstable modes, which is initially exponential, becomes only linear at the later times. And the abelianization works only in the exponential period of instability development. However, it might well be that the abelianization becomes more efficient when the dynamical effects beyond the Hard Loop approximation are taken into account .
### 2.5 Isotropization
When instabilites grow the systems becomes more isotropic because the Lorentz force acts on particle’s momenta and the growing fields generate an extra momentum.
To explain the mechanism I assume, as previously, that initially there is a momentum surplus in the $`z`$ direction. The fluctuating current tends to flow in the $`z`$ direction with the wave vector pointing in the $`x`$ direction. Since the magnetic field has a $`y`$ component, the Lorentz force, which acts on partons flying along the $`z`$ axis, pushes the partons in the $`x`$ direction where there is a momentum deficit.
The effect of isotropization due to the action of the Lorentz force is nicely seen in the classical simulation . In Fig. 6, which is taken from , there are shown diagonal components of the energy-momentum tensor
$$T^{\mu \nu }=\frac{d^3p}{(2\pi )^3}\frac{p^\mu p^\nu }{E_p}f(𝐩).$$
The initial momentum distribution is given by Eq. (9), and consequently $`T^{xx}=0`$ at $`t=0`$. As seen in Fig. 6, $`T^{xx}`$ exponentially grows.
The system isotropizes not only due to the effect of the Lorentz force but also due to the momentum carried by the growing field. As explained in detail in , the momentum of the field is oriented along the wave vector which points in the direction of the momentum deficit. This effect has not been numerically studied yet but it is clear that the effect is comparable to that of Lorentz force only for suffuciently large field amplitudes.
## 3 Azimuthal fluctuations
In the first part of my talk I have argued that the quark-gluon plasma becomes isotropic fast due to the magnetic instabilities. And it has been recently observed that the system with isotropic momentum distribution manifests a hydrodynamic collective behaviour. The question arises whether such an approximate hydrodynamics can be distinguished from the real hydrodynamics describing a system which is in a local thermodynamic equilibrium. In the second part of my talk I propose to address the question by studying the azimuthal fluctuations.
In relativistic heavy-ion collisions both at CERN SPS and BNL RHIC, one observes a sizable elliptic flow which is quantified by the second angular harmonics $`v_2`$ of the azimuthal distribution of final state hadrons . The phenomenon, which is sensitive to the collision early stage when the interaction zone is of the almond shape, is naturally explained within a hydrodynamics as a result of large density gradients . Hydrodynamic description requires that the system under study is in a local thermodynamical equilibrium. However, an approximate hydrodynamic behaviour occurs, as argued in , when the momentum distribution of liquid components is merely isotropic in the local rest frame. The point is that the structure of the ideal fluid energy-momentum tensor i.e. $`T^{\mu \nu }=(\epsilon +p)u^\mu u^\nu pg^{\mu \nu }`$, where $`\epsilon `$, $`p`$ and $`u^\mu `$ is the energy density, pressure and hydrodynamic velocity, respectively, holds for an arbitrary, though isotropic momentum distribution. $`\epsilon `$ and $`p`$ are then not the energy density and pressure but the moments of the distribution function which are equal the energy density and pressure in the equilibrium limit. Since the tensor $`T^{\mu \nu }`$ obeys the continuity equation $`_\mu T^{\mu \nu }=0`$, one gets an analogue of the Euler equation. However, due to the lack of thermodynamic equilibrium there is no entropy conservation and the equation of state is missing.
Usually, non-equilibrium fluctuations are significantly smaller than the equilibrium fluctuations of the same quantity. A specific example of such an situation has been discussed in Sec. 2.1. Therefore, I expect that the fluctuations of $`v_2`$ produced in the course of real hydrodynamic evolution are significantly smaller than those generated in the non-equilibrium quark-gluon plasma which is merely isotropic. It should be stressed here that the elliptic flow is generated in the collision early stage. Thus, I propose to carefully measure the fluctuations of $`v_2`$ as discussed in . Since such a measurement is rather difficult, I also consider an integral measurement of azimuthal fluctuations proposed in which can also help to distinguish the equilibrium from non-equilibrium fluctuations.
### 3.1 Elliptic flow fluctuations
In my discussion of $`v_2`$ fluctuations I follow where the standard method to measure the elliptic flow was used. The method focuses on the angular distributions relative to direction of the impact parameter. The experimental procedure splits in two steps which should be as independent as possible. In the first step, one determines the impact parameter direction $`\psi _R`$, while in the second step, one constructs the distribution of azimuthal angle relative to $`\psi _R`$ and one computes the Fourier coefficients.
The one-particle distribution in a single event can be written as
$$P_{\mathrm{ev}}(\varphi )=\frac{1}{2\pi }\left[1+2\underset{n=1}{\overset{\mathrm{}}{}}v_n\mathrm{cos}\left(n(\varphi \psi _R)\right)\right]\mathrm{\Theta }(\varphi )\mathrm{\Theta }(2\pi \varphi ).$$
(10)
Since the reaction plane is never reconstructed precisely and the real reaction plane angle $`\psi _R`$ deviates from the estimated angle $`\psi _E`$, the $`n`$th Fourier amplitude $`v_n`$ is determined as
$$v_n=\frac{1}{R_n}\overline{\mathrm{cos}\left(n(\varphi \psi _E)\right)},$$
where $`R_n\mathrm{cos}\left(n(\psi _R\psi _E)\right)`$ is the reaction plane resolution factor and $`\overline{\mathrm{}}`$ denotes averaging over particles from a single event.
Let me now consider an ensemble of events with every event representing a single nucleus-nucleus collision at fixed value (not direction) of the impact parameter. The angle $`\psi _R`$ (and $`\psi _E`$) obviously varies form event to event. The question is how to detect the event-by-event fluctuations of the Fourier amplitudes $`v_n`$. According to , the amplitude averaged over events is defined as
$$v_n\stackrel{\mathrm{def}}{=}\frac{\overline{\mathrm{cos}\left(\mathrm{n}(\varphi \psi _\mathrm{E})\right)}}{\mathrm{R}_\mathrm{n}},$$
where $`\mathrm{}`$ denotes averaging over events. We define the second moment as
$$v_n^2\stackrel{\mathrm{def}}{=}\frac{1}{\mathrm{R}_\mathrm{n}^2}\overline{\mathrm{cos}\left(\mathrm{n}(\varphi \psi _\mathrm{E})\right)}^2,$$
and the fluctuations are
$$\mathrm{𝖵𝖺𝗋}(v_n)v_n^2v_n^2=\frac{1}{R_n^2}\left(\overline{\mathrm{cos}\left(n(\varphi \psi _E)\right)}^2\overline{\mathrm{cos}\left(n(\varphi \psi _E)\right)}^2\right).$$
There are several sources of trivial $`v_2`$ fluctuations which are not related to the system’s dynamics of interest. I start with the fluctuations caused by a varying number of particles used to determine $`\psi _E`$ and $`v_2`$. I assume here that the Fourier amplitudes $`v_n`$ do not change from event to event and that the only correlations in the system are those due to the flow. Then, the azimuthal distribution of $`N`$ particles is a product of $`N`$ single particle distributions
$$P_{\mathrm{ev}}^N(\varphi _1,\varphi _2,\mathrm{},\varphi _N)=𝒫_NP_{\mathrm{ev}}(\varphi _1)P_{\mathrm{ev}}(\varphi _2)\mathrm{}P_{\mathrm{ev}}(\varphi _N),$$
(11)
where $`𝒫_N`$ is the multiplicity distribution while all distributions $`P_{\mathrm{ev}}(\varphi _i)`$ are given by Eq. (10). The single particle distributions $`P_{\mathrm{ev}}(\varphi _i)`$ are correlated to each other because of the common angle $`\psi _R`$.
Using the distribution (11), one finds (neglecting $`v_4`$) the variance of $`v_2`$ as
$$\mathrm{𝖵𝖺𝗋}(v_2)=\frac{1}{2R_2^2N}+v_2^2\frac{R_2^2R_2^2}{R_2^2},$$
(12)
which holds for $`N1`$ and small multiplicity fluctuations. The second term in r.h.s of Eq. (12) appears to be much smaller than the first one as $`R_2^2R_2^2M^2`$, where $`M`$ is the number of particles used to determine the reaction plane and $`M`$ is assumed to be of the same order as $`N`$. Thus, the statistical noise contribution to $`\delta v_2\sqrt{\mathrm{𝖵𝖺𝗋}(v_2)}`$ is finally estimated as
$$\delta v_2=\frac{1}{R_2\sqrt{2N}}.$$
(13)
As well known, $`v_2`$ strongly depends on the collision impact parameter $`b`$. Using the parameterisation of this dependence given in , one computes $`\delta v_2`$ as $`(dv_2/db)\delta b`$. The impact parameter is measurable through the multiplicity $`N_p`$ of participating nucleons. $`N_p`$ is directly related to $`b`$. For $`b10`$ fm, when the flow in Au-Au collisions is maximal, the $`v_2`$ fluctuations due to the impact parameter variation vanish because $`dv_2/db=0`$. For $`b=5`$ fm, where $`v_20.03`$, one finds $`\delta v_2810^4\delta N_p`$. When $`\delta N_p=30`$ and $`N=500`$, the magnitude of the $`v_2`$ fluctuations caused by the impact parameter variation is approximately equal to that of the statistical noise.
The next source of trivial $`v_2`$ fluctuations is a variation of thermodynamic parameters. The contribution caused by the particle number fluctuations can be estimated as
$$\delta v_2=\frac{dv_2}{dN}\delta N=v_2\frac{\delta N}{N}P,$$
where the effective power $`P`$ is
$$P\frac{d\mathrm{ln}v_2}{d\mathrm{ln}N}=\frac{N}{v_2}\frac{dv_2}{dN}.$$
Assuming the poissonian character of multiplicity fluctuations, one obtains
$$\delta v_2=\frac{v_2}{\sqrt{N}}P.$$
(14)
The value of the index $`P`$ can be estimated within the hydrodynamic model as $`P0.4`$. Comparing Eqs. (13) and (14), one finds that the ratio of the thermodynamic fluctuations to the statistical noise, is 0.04 for $`v_2=0.07`$ and $`P=0.4`$. Thus, the thermodynamic fluctuations are much smaller than the statistical noise.
Concluding this section, I propose to perform a systematic measurement of event-by-event $`v_2`$ fluctuations for the centrality corresponding to the maximal elliptic flow. Then, the fluctuations caused by the impact parameter variation vanish or at least they are very small. If the flow is built up in the course of hydrodynamic evolution of the equilibrium system, $`v_2`$ should be dominated by the statistical noise related to the finite particle number. The noise can be identified due to the characteristic $`1/N`$ dependence.
### 3.2 $`\mathrm{\Phi }`$ measure of azimuthal fluctuations
Since a measurement of $`v_2`$ fluctuations discussed in the previous section is rather difficult, I suggest to consider a much simpler integral measurement of azimuthal fluctuations , using the so-called $`\mathrm{\Phi }`$ measure introduced in .
The correlation (or fluctuation) measure $`\mathrm{\Phi }`$ is defined as follows. One defines the variable $`z\stackrel{\mathrm{def}}{=}\mathrm{x}\overline{\mathrm{x}}`$, where $`x`$ is a single particle characteristics such as the particle transverse momentum or the azimuthal angle. In this section the overline does not denote averaging over particles from a single event but averaging over a single particle inclusive distribution. $`x`$ is identified here with the particle azimuthal angle. The event variable $`Z`$, which is a multiparticle analog of $`z`$, is defined as $`Z\stackrel{\mathrm{def}}{=}_{\mathrm{i}=1}^\mathrm{N}(\mathrm{x}_\mathrm{i}\overline{\mathrm{x}})`$, where the summation runs over particles from a given event. By construction, $`Z=0`$. The measure $`\mathrm{\Phi }`$ is finally defined as
$$\mathrm{\Phi }\stackrel{\mathrm{def}}{=}\sqrt{\frac{\mathrm{Z}^2}{\mathrm{N}}}\sqrt{\overline{\mathrm{z}^2}}.$$
$`\mathrm{\Phi }`$ obviously vanishes in the absence of any inter-particle correlations. Other properties of $`\mathrm{\Phi }`$ are discussed in .
The $`\mathrm{\Phi }`$ measure is sensitive to the azimthal fluctuations caused by the transverse collective flow. Let me compute it, assuming that the only correlations present in the system are due to the collective flow. The inclusive $`\varphi `$ distrubtion, which is flat in the range $`[0,2\pi ]`$, provides $`\overline{\varphi }=\pi `$ and $`\overline{\varphi ^2}=\frac{4}{3}\pi ^2`$, and consequently, $`\overline{z^2}=\frac{1}{3}\pi ^2`$. Since $`Z=_{i=1}^N(\varphi _i\overline{\varphi })`$, one computes $`Z^2`$, using the event distribution (10), as
$`Z^2={\displaystyle _0^{2\pi }}{\displaystyle \frac{d\psi _R}{2\pi }}{\displaystyle \underset{N}{}}𝒫_N{\displaystyle _0^{2\pi }}𝑑\varphi _1\mathrm{}{\displaystyle _0^{2\pi }}𝑑\varphi _NP_{\mathrm{ev}}(\varphi _1)\mathrm{}P_{\mathrm{ev}}(\varphi _N)(\varphi _1+\mathrm{}+\varphi _NN\overline{\varphi })^2.`$ (15)
The averaging over the amplitudes $`v_n`$, which is not shown here, is implied. At first glance, the multi-particle distribution from Eq. (15) might look as a simple product of one-particle distributions. One should note however that every $`P_{\mathrm{ev}}(\varphi )`$ depends on the reaction plane angle $`\psi _R`$, and the integration over $`\psi _R`$ leads to the correlated multi-particle distribution.
After elementary calculation, one finds for small $`v_n`$, poissonian multiplicity distribution and $`N1`$, an approximate expression of interest
$$\mathrm{\Phi }\frac{3}{\pi ^2}N\underset{n=1}{\overset{\mathrm{}}{}}\left(\frac{v_n}{n}\right)^2.$$
(16)
If all amplitudes $`v_n`$ except $`v_2`$ vanish, as it approximately happens in the central rapidity domain, and $`v_2`$ equals a unique value 0.07, Eq. (16) for $`N=500`$ gives $`\mathrm{\Phi }=0.2`$.
As already mentioned, the transverse flow is far not the only source of the azimuthal fluctuations. In particular, the effect of quantum statistics contribute here. To estimate the effect one computes $`\mathrm{\Phi }`$ in the equilibrium ideal quantum gas. The result reads
$$\mathrm{\Phi }=\frac{\pi }{\sqrt{3}}\left(\sqrt{\frac{\stackrel{~}{\rho }}{\rho }}1\right),$$
(17)
where
$$\rho \frac{d^3p}{(2\pi )^3}\frac{1}{\lambda ^1e^{\beta E_p}\pm 1},\stackrel{~}{\rho }\frac{d^3p}{(2\pi )^3}\frac{\lambda ^1e^{\beta E_p}}{(\lambda ^1e^{\beta E_p}\pm 1)^2},$$
$`\lambda `$ denotes the fugacity, the upper sign is for fermions and the lower one for bosons. As seen, $`\mathrm{\Phi }`$ is independent of the system’s volume and of the number of particle’s internal degrees of freedom. For massless bosons with vanishing chemical potential, Eq. (17) gives $`\mathrm{\Phi }0.3`$ for any $`T`$. More realistic calculations provide $`\mathrm{\Phi }0.06`$ for chemically equilibrated pions at $`T=150`$ MeV.
I conclude this section by saying that a measured value of $`\mathrm{\Phi }`$, which would significantly exceed predictions of Eq. (16) with non-fluctuating amplitudes $`v_n`$, would be an obvious signal of sizable dynamical fluctuations.
## 4 Final remarks
The magnetic instabilities provide a plausible mechanism responsible for a surprisingly short equilibration time observed in relativistic heavy-ion collisions. Fast isotropization is a distinctive feature of the mechanism. It is certainly desirable to look for experimentally detectable signals of the instabilities driven thermalization. In my talk I have proposed to study azimuthal fluctuations, in particular the event-by-event fluctuations of the elliptic flow which is generated at the collision early stage. I have not been able to present a quantitative prediction but observation of sizeable dynamical fluctuations would be a strong argument that behind a smooth hydrodynamic evolution there is a violent phenomenon of plasma instabilities.
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# Thermodynamics of Asymptotically Flat Charged Black Holes in Third Order Lovelock Gravity
## I Introduction
Both string theory as well as brane world cosmology assume that spacetime possesses more than four dimensions. In string theory, extra dimensions were promoted from an interesting curiosity to a theoretical necessity since superstring theory requires a ten-dimensional spacetime to be consistent from the quantum point of view, while in the brane world cosmology (which is also consistent with string theory) the matter and gauge interaction are localized on a 3-brane, embedded into a higher dimensional spacetime and gravity propagates in the whole of spacetime. This underscores the need to consider gravity in higher dimensions. In this context one may use another consistent theory of gravity in any dimension with a more general action. This action may be written, for example, through the use of string theory. The effect of string theory on classical gravitational physics is usually investigated by means of a low energy effective action which describes gravity at the classical level Wit1 . This effective action consists of the Einstein-Hilbert action plus curvature-squared terms and higher powers as well, and in general give rise to fourth order field equations and bring in ghosts. However, if the effective action contains the higher powers of curvature in particular combinations, then only second order field equations are produced and consequently no ghosts arise Zw . The effective action obtained by this argument is precisely of the form proposed by Lovelock Lov :
$$I_G=d^dx\sqrt{g}\underset{k=0}{\overset{[d/2]}{}}\alpha _k_k$$
(1)
where $`[z]`$ denotes integer part of $`z`$, $`\alpha _k`$ is an arbitrary constant and $`_k`$ is the Euler density of a $`2k`$-dimensional manifold,
$$_k=\frac{1}{2^k}\delta _{\rho _1\sigma _1\mathrm{}\rho _k\sigma _k}^{\mu _1\nu _1\mathrm{}\mu _k\nu _k}R_{\mu _1\nu _1}^{\rho _1\sigma _1}\mathrm{}R_{\mu _k\nu _k}^{\rho _k\sigma _k}$$
(2)
In Eq. (2) $`\delta _{\rho _1\sigma _1\mathrm{}\rho _k\sigma _k}^{\mu _1\nu _1\mathrm{}\mu _k\nu _k}`$ is the generalized totally anti-symmetric Kronecker delta and $`R_{\mu \nu }^{\rho \sigma }`$ is the Riemann tensor. It is worthwhile to mention that in $`d`$ dimensions, all terms for which $`k>[d/2]`$ are identically equal to zero, and the term $`k=d/2`$ is a topological term. So, only terms for which $`k<d/2`$ are contributing to the field equations.
In this paper we want to restrict ourself to the first four terms of Lovelock gravity. The first term is the cosmological term which we ignore it in the investigation of the properties and thermodynamics of the solutions. The second term is the Einstein term, and the third and fourth terms are the second order Lovelock (Gauss-Bonnet) and third order Lovelock terms respectively. From a geometric point of view, the combination of these terms in seven-dimensional spacetimes, is the most general Lagrangian producing second order field equations, as in the four-dimensional gravity which the Einstein-Hilbert action is the most general Lagrangian producing second order field equations. Here, we will obtain asymptotically flat solutions of third order Lovelock gravity and investigate their thermodynamics.
Indeed, it is interesting to explore black holes in generalized gravity theories in order to discover which properties are peculiar to Einstein’s gravity, and which are robust features of all generally covariant theories of gravity. Due to the nonlinearity of the field equations, it is very difficult to find out nontrivial exact analytical solutions of Einstein’s equation with the higher curvature terms. In most cases, one has to adopt some approximation methods or find solutions numerically. However, exact static spherically symmetric black hole solutions of the Gauss-Bonnet gravity have been found in Ref. Des ; Whe , and of the Einstein-Maxwell-Gauss-Bonnet and Einstein-Born-Infeld-Gauss-Bonnet models in Refs. Wil1 ; Wil2 . Black hole solutions with nontrivial topology in this theory have been also studied in Refs. Cai ; Ish . The thermodynamics of the uncharged static spherically black hole solutions has been considered in MS and of charged solutions in Wil2 ; Odin . All of these known solutions in Gauss-Bonnet gravity are static. Recently one of us has introduced two new classes of rotating solutions of second order Lovelock gravity and investigate their thermodynamics Deh1 .
Also the static spherically symmetric solutions of the dimensionality continued gravity have been explored in Ref. Zan , while black hole solutions with nontrivial topology in this theory have been studied in Ref. Aros . The thermodynamics of these solutions have been investigated in Refs. Mun ; Chr ; Clu . Up to our knowledge, no asymptotically flat solution for Lovelock gravity higher than second order (Gauss-Bonnet) has been obtained till now with two or more fundamental constants. Indeed, the asymptotically flat solution of Ref. Chr has only one fundamental constant which is the gravitational constant $`G`$. In this paper we want to find new static solutions of third order Lovelock gravity which are asymptotically flat and contain two and three fundamental constants and investigate their thermodynamics. As we will show later, these asymptotically flat solutions have some properties which do not occur in Einstein or Gauss-Bonnet gravity.
The outline of our paper is as follows. We give a brief review of the field equations of third order Lovelock gravity in Sec. II. In Sec. III, we present the static solutions of third order Lovelock gravity in the presence of electromagnetic field with special values of $`\alpha _2`$ and $`\alpha _3`$, and investigate their properties. In Sec. IV we obtain mass, entropy, temperature, charge, and electric potential of the $`d`$-dimensional black hole solutions and show that these quantities satisfy the first law of thermodynamics. We also perform a local stability analysis of the black holes in the canonical and grand canonical ensembles. In Sec. V, we introduce the general asymptotically flat solutions with three fundamental constants and investigate their thermodynamics. We finish our paper with some concluding remarks.
## II Field equations
The main fundamental assumptions in standard general relativity are the requirements of general covariance and that the field equations for the metric be second order. Based on the same principles, the Lovelock Lagrangian is the most general Lagrangian in classical gravity which produces second order field equations for the metric. The action of third order Lovelock gravity in the presence of electromagnetic field may be written as
$$I=d^dx\sqrt{g}\left(2\mathrm{\Lambda }+_1+\alpha _2_2+\alpha _3_3F_{\mu \nu }F^{\mu \nu }\right)$$
(3)
where $`\mathrm{\Lambda }`$ is the cosmological constant, $`\alpha _2`$ and $`\alpha _3`$ are Gauss-Bonnet and third order lovelock coefficients, $`_1=R`$ is just the Einstein-Hilbert Lagrangian, $`_2=R_{\mu \nu \gamma \delta }R^{\mu \nu \gamma \delta }4R_{\mu \nu }R^{\mu \nu }+R^2`$ is the Gauss-Bonnet Lagrangian, and
$`_3`$ $`=`$ $`2R^{\mu \nu \sigma \kappa }R_{\sigma \kappa \rho \tau }R_{\mu \nu }^{\rho \tau }+8R_{\sigma \rho }^{\mu \nu }R_{\nu \tau }^{\sigma \kappa }R_{\mu \kappa }^{\rho \tau }+24R^{\mu \nu \sigma \kappa }R_{\sigma \kappa \nu \rho }R_\mu ^\rho `$ (4)
$`+3RR^{\mu \nu \sigma \kappa }R_{\sigma \kappa \mu \nu }+24R^{\mu \nu \sigma \kappa }R_{\sigma \mu }R_{\kappa \nu }+16R^{\mu \nu }R_{\nu \sigma }R_\mu ^\sigma 12RR^{\mu \nu }R_{\mu \nu }+R^3`$
is the third order Lovelock Lagrangian. In Eq. (3) $`F_{\mu \nu }=_\mu A_\nu _\nu A_\mu `$ is electromagnetic tensor field and $`A_\mu `$ is the vector potential.
Since in Lovelock gravity, only terms for which $`k<d/2`$ are contributing to the field equations and we want to consider the third order lovelock gravity, therefore we consider the $`d`$-dimensional spacetimes with $`d7`$. Varying the action with respect to the metric tensor $`g_{\mu \nu }`$ and electromagnetic tensor field $`F_{\mu \nu }`$ the equations of gravitation and electromagnetic fields are obtained as:
$`G_{\mu \nu }^{(1)}+\mathrm{\Lambda }g_{\mu \nu }+\alpha _2G_{\mu \nu }^{(2)}+\alpha _3G_{\mu \nu }^{(3)}=T_{\mu \nu }`$ (5)
$`_\nu F^{\mu \nu }=0`$ (6)
where $`T_{\mu \nu }=2F_\mu ^\rho F_{\rho \nu }\frac{1}{2}F_{\rho \sigma }F^{\rho \sigma }g_{\mu \nu }`$ is the energy-momentum tensor of electromagnetic field, $`G_{\mu \nu }^{(1)}`$ is just the Einstein tensor, and $`G_{\mu \nu }^{(2)}`$ and $`G_{\mu \nu }^{(3)}`$ are given as Hoi :
$`G_{\mu \nu }^{(2)}`$ $`=`$ $`2(R_{\mu \sigma \kappa \tau }R_\nu ^{\kappa \tau \sigma }2R_{\mu \rho \nu \sigma }R^{\rho \sigma }2R_{\mu \sigma }R_\nu ^\sigma +RR_{\mu \nu }){\displaystyle \frac{1}{2}}_2g_{\mu \nu }`$
$`G_{\mu \nu }^{(3)}`$ $`=`$ $`3(4R^{\tau \rho \sigma \kappa }R_{\sigma \kappa \lambda \rho }R_{\nu \tau \mu }^\lambda 8R_{\lambda \sigma }^{\tau \rho }R_{\tau \mu }^{\sigma \kappa }R_{\nu \rho \kappa }^\lambda +2R_\nu ^{\tau \sigma \kappa }R_{\sigma \kappa \lambda \rho }R_{\tau \mu }^{\lambda \rho }`$
$`R^{\tau \rho \sigma \kappa }R_{\sigma \kappa \tau \rho }R_{\nu \mu }+8R_{\nu \sigma \rho }^\tau R_{\tau \mu }^{\sigma \kappa }R_\kappa ^\rho +8R_{\nu \tau \kappa }^\sigma R_{\sigma \mu }^{\tau \rho }R_\rho ^\kappa `$
$`+4R_\nu ^{\tau \sigma \kappa }R_{\sigma \kappa \mu \rho }R_\tau ^\rho 4R_\nu ^{\tau \sigma \kappa }R_{\sigma \kappa \tau \rho }R_\mu ^\rho +4R^{\tau \rho \sigma \kappa }R_{\sigma \kappa \tau \mu }R_{\nu \rho }+2RR_\nu ^{\kappa \tau \rho }R_{\tau \rho \kappa \mu }`$
$`+8R_{\nu \mu \rho }^\tau R_\sigma ^\rho R_\tau ^\sigma 8R_{\nu \tau \rho }^\sigma R_\sigma ^\tau R_\mu ^\rho 8R_{\sigma \mu }^{\tau \rho }R_\tau ^\sigma R_{\nu \rho }4RR_{\nu \mu \rho }^\tau R_\tau ^\rho `$
$`+4R^{\tau \rho }R_{\rho \tau }R_{\nu \mu }8R_\nu ^\tau R_{\tau \rho }R_\mu ^\rho +4RR_{\nu \rho }R_\mu ^\rho R^2R_{\nu \mu }){\displaystyle \frac{1}{2}}_3g_{\mu \nu }`$
## III Static Solutions
The metric of $`d`$-dimensional static spherically symmetric spacetime and the vector potential may be written as:
$`ds^2`$ $`=`$ $`f(r)dt^2+{\displaystyle \frac{dr^2}{f(r)}}+r^2d\mathrm{\Omega }^2`$
$`A_\mu `$ $`=`$ $`h(r)\delta _\mu ^t`$ (7)
where $`d\mathrm{\Omega }^2`$ is the metric of a $`(d2)`$-sphere. The functions $`f(r)`$ and $`h(r)`$ may be obtained by solving the field equations (5) and (6).
### III.1 Seven-dimensional Solutions
As stated before, the gravitational field equation of third order Lovelock gravity in seven dimensions is the most general second order differential equation which presents the solutions of gravity. Indeed, the solution of third order Lovelock gravity in seven dimensions is the most general solution of gravity, based on the principle of standard general relativity. Therefore, first we obtain the seven-dimensional solutions of third order Lovelock gravity in the presence of electromagnetic field and investigate their properties. Using Eq. (6) one can show that $`h(r)=q/(4r^4)`$, where $`q`$ is an arbitrary real constant which is related to the charge of the solution. To find the function $`f(r)`$, one may use any components of Eq. (5). The simplest equation is the $`rr`$-component of these equations which can be written as
$$\left(3\alpha _3(f1)^2+\alpha _2(f1)r^2\frac{r^4}{24}\right)r^5f^{}+\alpha _2r^6(f1)^2\frac{r^8}{6}(f1)\frac{\mathrm{\Lambda }}{60}r^{10}=\frac{q^2}{60}$$
(8)
where prime denotes the derivative with respect to $`r`$. We consider the solutions of Eq. (8) for any arbitrary values of $`\alpha _2`$ and $`\alpha _3`$ in Sec. V. Here, we study the special case of $`\alpha _3=2\alpha _{2}^{}{}_{}{}^{2}=\alpha ^2/72`$. Even for this special case, there exist three fundamental constants, $`\mathrm{\Lambda }`$, $`G=1`$ and $`\alpha `$ in the solution, while the solution of ref. Chr has two fundamental constants. As we will see below this solution has some properties which will not happen in the Gauss-Bonnet gravity with three fundamental constants. The real solution of Eq. (8) with the above values of $`\alpha _2`$ and $`\alpha _3`$ is
$$f(r)=1+\frac{r^2}{\alpha }\left\{1\left(1+\frac{\mathrm{\Lambda }\alpha }{5}+\frac{3\alpha m}{5r^6}\frac{3\alpha q^2}{10r^{10}}\right)^{1/3}\right\}$$
(9)
In the above equations $`m`$ is an integration constant which is related to the mass of the solution. Unlike the solutions in Gauss-Bonnet gravity which have two branches, here the solution (9) has only one branch. Indeed, Eq. (8) with the above $`\alpha `$’s has the real solution (9) and two complex solutions which are the complex conjugate of each other. This feature is the same as the asymptotically anti-de Sitter (AdS) solution of Ref. Chr which has a unique anti de Sitter vacuum.
In order to study the general structure of this solution, we first look for the asymptotic behavior of the solutions. It is easy to find out that these solutions are asymptotically flat for $`\mathrm{\Lambda }=0`$, AdS for $`\mathrm{\Lambda }<0`$ and de Sitter (dS) for $`\mathrm{\Lambda }>0`$. In this paper we are interested in the case of asymptotically flat solutions, and therefore we put $`\mathrm{\Lambda }=0`$. One can show that the Kretschmann scalar $`R_{\mu \nu \lambda \kappa }R^{\mu \nu \lambda \kappa }`$ diverges at $`r=0`$, and therefore there is a curvature singularity located at $`r=0`$. Now we look for the existence of horizons, and therefore we look for possible black hole solutions. The horizons, if any exist, are given by the zeros of the function $`f(r)=g^{rr}`$.
First, we consider the uncharged solutions, where the horizon(s) is located at
$`r_+`$ $`=`$ $`\left\{{\displaystyle \frac{\alpha }{2}}+{\displaystyle \frac{\sqrt{180m75\alpha ^2}}{30}}\right\}^{1/2}`$ (10)
$`r_{}`$ $`=`$ $`\left\{{\displaystyle \frac{\alpha }{2}}{\displaystyle \frac{\sqrt{180m75\alpha ^2}}{30}}\right\}^{1/2}`$ (11)
As one can see from Eqs. (10), and (11), for $`\alpha >0`$ there exists only one horizon provided the mass parameter, $`m`$, is greater than $`m_{\mathrm{crit}}=5/3\alpha ^2`$. That is, there exist a minimum value for the mass (or radius of horizon) in order to have uncharged asymptotically flat black hole solution. More interesting is the uncharged solutions of third order Lovelock gravity with $`\alpha <0`$. In this case there exist an extreme value for the mass parameter, $`m_{\mathrm{ext}}=5/12\alpha ^2`$. The uncharged solution presents a black hole with inner and outer horizon provided $`m>m_{\mathrm{ext}}`$, an extreme black hole for $`m=m_{\mathrm{ext}}`$ and a naked singularity otherwise. This property happens only for third order Lovelock gravity and does not occur in Einstein or Gauss-Bonnet gravity. Thus, it is natural to suppose that new solutions in higher order Lovelock gravity might provide us with a new window on some corners of low energy limit of string theory.
Second, we consider the charged solutions. It is easy to show that the solution presents a black hole solution with two inner and outer horizon, provided the charge parameter $`q`$ is less than, $`q_{\mathrm{ext}}`$, an extreme black hole for $`q=q_{\mathrm{ext}}`$ and a naked singularity otherwise, where $`q_{\mathrm{ext}}`$ is
$$q_{\mathrm{ext}}=\frac{1}{240}\left\{5760m^2+13200\alpha ^2m7625\alpha ^4+5\alpha \sqrt{15(96m25\alpha ^2)^3}\right\}$$
(12)
### III.2 $`d`$-dimensional Solutions
The $`d`$-dimensional static solutions in third order Lovelock gravity may be obtained by solving Eqs. (5) and (6) for the metric given in Eq. (7). Using Eq. (6), one can show that $`h(r)=q/[(d3)r^{d3}]`$, where $`q`$ is an arbitrary real constant which is related to the charge of the solution, and the function $`f(r)`$ is the solution of the following equation:
$`[180(_5^{d2})\alpha _3r(f1)^26(_3^{d2})\alpha _2(f1)r^3+{\displaystyle \frac{d2}{2}}r^5]f^{}+\mathrm{\Lambda }r^6`$
$`+360(_6^{d2})\alpha _3(f1)^312(_4^{d2})\alpha _2r^2(f1)^2+(_2^{d2})r^4(f1)=q^2r^{102d}`$ (13)
where prime denotes the derivative with respect to $`r`$. The solution of Eq. (13) for arbitrary values of $`\alpha _2`$ and $`\alpha _3`$ will be introduced in Sec. V. Here we consider the solutions of Eq. (13) for $`\alpha _i`$’s given as
$$\alpha _3=\frac{\alpha ^2}{72(_4^{d3})},\alpha _2=\frac{\alpha }{(d3)(d4)}$$
(14)
Again, Eq. (13) with condition (14) has one real and two complex solutions which are the complex conjugate of each other. The real solution of Eq. (13) with condition (14) is
$$f(r)=1+\frac{r^2}{\alpha }\left\{1\left(1+\frac{6\mathrm{\Lambda }\alpha }{(d1)(d2)}+\frac{3\alpha m}{(d2)r^{d1}}\frac{6\alpha q^2}{\left(d2\right)(d3)r^{2d4}}\right)^{1/3}\right\}$$
(15)
where $`m`$ is the mass parameter.
These solutions are asymptotically flat for $`\mathrm{\Lambda }=0`$ and AdS or dS for negative or positive values of $`\mathrm{\Lambda }`$ respectively. Also one can show that the Kretschmann scalar diverges at $`r=0`$ and therefore there is a curvature singularity located at $`r=0`$.
In order to investigate the existence of black hole solutions for $`\mathrm{\Lambda }=0`$ we, first, consider the uncharged solutions. As in the case of seven-dimensional solutions, one can show that for $`\alpha >0`$, there exists a critical value of horizon radius, $`r_{+\mathrm{crit}}`$ (or a critical value of mass, $`m_{\mathrm{crit}}`$) such that the solution presents a black hole for $`r_+>r_{+\mathrm{crit}}`$, and a naked singularity otherwise, where $`r_{+\mathrm{crit}}`$ is the larger real root of the following equation:
$$3(d3)r_+^4+3(d5)\alpha r_+^2+(d7)\alpha ^2=0$$
(16)
The critical mass, $`m_{\mathrm{crit}}`$ now may be obtained easily as
$`m_{\mathrm{crit}}`$ $`=`$ $`{\displaystyle \frac{(d2)\left[d+3+\sqrt{3(d1)(d9)}\right]}{3}}`$ (17)
$`\times \left\{{\displaystyle \frac{3d15\sqrt{3(d1)(d9)}}{6}}\right\}^{(d7)/2}\left({\displaystyle \frac{\alpha }{d3}}\right)^{(d3)/2}`$
Note that the critical mass is real only for $`d9`$, and therefore there exist no uncharged black hole in dimension greater than nine for $`\alpha _2`$ and $`\alpha _3`$ given in Eqs. (14). Unlike to this special case, we will see in Sec. V that black hole solutions exist in any dimension for the arbitrary values of $`\alpha _2`$ and $`\alpha _3`$.
For $`\alpha <0`$, the $`d`$-dimensional solutions of third order Lovelock gravity present a new property that does not occur in lower order of Lovelock gravity. Indeed the uncharged solution for negative values of $`\alpha `$ presents a black hole with two inner and outer horizon, provided $`r_+>r_{+\mathrm{ext}}`$, an extreme black hole for $`r_+=r_{+\mathrm{ext}}`$, and a naked singularity otherwise, where $`r_{+\mathrm{ext}}`$ is the smaller solution of Eq. (17). In this case the extremal value of mass is
$`m_{\mathrm{ext}}`$ $`=`$ $`{\displaystyle \frac{(d2)\left[d+3\sqrt{3(d1)(d9)}\right]}{3}}`$ (18)
$`\left\{{\displaystyle \frac{3d15+\sqrt{3(d1)(d9)}}{6}}\right\}^{(d7)/2}\left({\displaystyle \frac{\left|\alpha \right|}{d3}}\right)^{(d3)/2}`$
Second, we consider the charged solutions. It is easy to show that the solution presents a black hole solution with two inner and outer horizon, provided $`q`$ is less than $`q_{\mathrm{ext}}`$, an extreme black hole for $`q=q_{\mathrm{ext}}`$, and a naked singularity otherwise , where $`q_{\mathrm{ext}}`$ is the solution of
$$3(d5)\alpha r_+^{2d8}+3(d3)r_+^{2d6}+(d7)\alpha ^2r_+^{2d10}6(d2)^1q^2=0$$
(19)
## IV Thermodynamics of black holes
One can obtain the temperature of the event horizon by analytic continuation of the metric. The analytical continuation of the Lorentzian metric by $`ti\tau `$ yields the Euclidean section, whose regularity at $`r=r_+`$ requires that we should identify $`\tau \tau +\beta _+`$, where $`\beta _+`$ is the inverse Hawking temperature of the horizon given as
$$\beta _+^1=T_+=\frac{f^{}(r)}{4\pi }=\frac{(d2)r_+^{2d10}[3(d3)r_+^4+3(d5)\alpha r_+^2+(d7)\alpha ^2]2q^2}{12\pi (r_+^2+\alpha )^2r_+^{2d9}}$$
(20)
Usually entropy of the black holes satisfies the so-called area law of entropy which states that the black hole entropy equals to one-quarter of horizon area Beck . One of the surprising and impressive feature of this area law of entropy is its universality. It applies to all kind of black holes and black strings of Einstein gravity Haw . However, in higher derivative gravity the area law of entropy is not satisfied in general fails . It is known that the entropy in Lovelock gravity is Myer
$$S=\frac{1}{4}\underset{k=1}{\overset{[(d1)/2]}{}}k\alpha _kd^{d2}x\sqrt{\stackrel{~}{g}}\stackrel{~}{}_{k1}$$
(21)
where the integration is done on the $`(d2)`$-dimensional spacelike hypersurface of Killing horizon, $`\stackrel{~}{g}_{\mu \nu }`$ is the induced metric on it, $`\stackrel{~}{g}`$ is the determinant of $`\stackrel{~}{g}_{\mu \nu }`$ and $`\stackrel{~}{}_k`$ is the $`k`$th order Lovelock Lagrangian of $`\stackrel{~}{g}_{\mu \nu }`$. Thus, the entropy in third order Lovelock gravity is
$$S=\frac{1}{4}d^{d2}x\sqrt{\stackrel{~}{g}}\left(1+2\alpha _2\stackrel{~}{R}+3\alpha _3(\stackrel{~}{R}_{\mu \nu \sigma \kappa }\stackrel{~}{R}^{\mu \nu \sigma \kappa }\stackrel{~}{R}_{\mu \nu }\stackrel{~}{R}^{\mu \nu }+\stackrel{~}{R}^2)\right)$$
(22)
where $`\stackrel{~}{R}_{\mu \nu \rho \sigma }`$ and $`\stackrel{~}{R}_{\mu \nu }`$ are Riemann and Ricci tensors and $`\stackrel{~}{R}`$ is the Ricci scalar for the induced metric $`\stackrel{~}{g}_{ab}`$ on the $`(d2)`$-dimensional horizon. It is a matter of calculation to show that the entropy of black holes is
$$S=\frac{1}{4}r_+^{d2}\mathrm{\Sigma }_{d2}\left(1+\frac{4\alpha _2(_2^{d2})}{r_+^2}+\frac{72(_4^{d2})\alpha _3}{r_+^4}\right)$$
(23)
where $`\mathrm{\Sigma }_{d2}`$ is the area of a unit radius $`(d2)`$-dimensional sphere.
The charge of the black hole can be found by calculating the flux of the electric field at infinity, yielding
$$Q=\frac{\mathrm{\Sigma }_{d2}}{4\pi }q$$
(24)
The electric potential $`\mathrm{\Phi }`$, measured at infinity with respect to the horizon, is defined by Cal
$$\mathrm{\Phi }=A_\mu \chi ^\mu |{}_{r\mathrm{}}{}^{}A_\mu \chi ^\mu |_{r=r_+},$$
(25)
where $`\chi =/t`$ is the null generator of the horizon. One finds
$$\mathrm{\Phi }=\frac{q}{(d3)r_+^{d3}}$$
(26)
The mass of black hole can be obtained by using the behavior of the metric at large $`r`$. It is easy to show that the mass of black hole is
$$M=\frac{\mathrm{\Sigma }_{d2}}{16\pi }m=\frac{(d2)(d3)(3r_+^4+3\alpha r_+^2+\alpha ^2)r_+^{2d10}+2q^2}{(d3)r_+^{d3}}$$
(27)
We now investigate the first law of thermodynamics. Using the expression for the entropy, the charge and the mass given in Eqs. (23), (24) and (27), one can compute $`M/r`$, $`S/r`$ and $`Q/r`$. Then, by using the chain rule, it is easy to show that the quantities
$$T=\left(\frac{M}{S}\right)_Q,\mathrm{\Phi }=\left(\frac{M}{Q}\right)_S$$
(28)
are exactly the same as the temperature and electric potential given in Eqs. (20) and (26) respectively. Thus, the thermodynamic quantities calculated in Eqs. (20) and (26) satisfy the first law of thermodynamics,
$$dM=TdS+\mathrm{\Phi }dQ$$
(29)
### IV.1 Stability in the canonical and the grand-canonical ensemble
The stability of a thermodynamic system with respect to the small variations of the thermodynamic coordinates, is usually performed by analyzing the behavior of the entropy $`S(M,Q)`$ around the equilibrium. The local stability in any ensemble requires that $`S(M,Q)`$ be a convex function of their extensive variables or its Legendre transformation must be a concave function of their intensive variables. Thus, the local stability can in principle be carried out by finding the determinant of the Hessian matrix of $`S`$ with respect to its extensive variables, $`𝐇_{X_iX_j}^S=[^2S/X_iX_j]`$, where $`X_i`$’s are the thermodynamic variables of the system. Indeed, the system is locally stable if the determinant of Hessian matrix satisfies $`𝐇_{X_i,X_j}^S0`$ Cvet . Also, one can perform the stability analysis through the use of the determinant of Hessian matrix of the energy with respect to its thermodynamic variables, and the stability requirement $`𝐇_{X_i,X_j}^S0`$ may be rephrased as $`𝐇_{Y_i,Y_j}^M0`$ Gub .
The number of the thermodynamic variables depends on the ensemble which is used. In the canonical ensemble, the charge is a fixed parameter, and therefore the positivity of the heat capacity $`C_Q=T(S/T)_Q`$ is sufficient to assure the local stability. The heat capacity for the solutions given by Eq. (15) shows that there exist an upper limit for the radius of horizon, $`r_{+\mathrm{st}}`$ which is the real root of $`St_{\mathrm{can}}=0`$, where $`St_{\mathrm{can}}`$ is
$`St_{\mathrm{can}}`$ $`=`$ $`(d2)[(d7)\alpha ^3+2r_{+\mathrm{st}}^{}{}_{}{}^{2}(d10)\alpha ^218r_{+\mathrm{st}}^{}{}_{}{}^{4}\alpha +3r_{+\mathrm{st}}^{}{}_{}{}^{6}(d3)]`$ (30)
$`6(2d5)r_{+\mathrm{st}}^{}{}_{}{}^{2d+12}q^26(2d9)r_{+\mathrm{st}}^{}{}_{}{}^{2d+10}\alpha q^2`$
Therefore, the charged solutions are stable provided the horizon radius lies in the range $`r_{+\mathrm{ext}}<r_+<r_{+\mathrm{st}}`$, where $`r_{+\mathrm{ext}}`$ is the solution of Eq. (19), and $`r_{+\mathrm{st}}`$ is the real solution of Eq. (30). Indeed, for $`r_+<r_{+\mathrm{ext}}`$, we have no black hole while for $`r_+>r_{+\mathrm{st}}`$, the black hole is not stable and therefore we have only an intermediate stable phase. Note that for the case of uncharged black holes and $`\alpha >0`$, $`r_+`$ should be greater than $`r_{+crit}`$ as stated before. Thus for positive $`\alpha `$, the horizon radius of stable black holes lies in the range $`r_{+\mathrm{crit}}<r_+<r_{+\mathrm{st}}`$, where $`r_{+\mathrm{crit}}`$ is the larger root of Eq. (17). Of course, for negative $`\alpha `$, which one can have a black hole with inner and outer horizon, the black hole is stable provided $`r_+`$ lies in the range $`r_{+\mathrm{ext}}<r_+<r_{+\mathrm{st}}`$, where now $`r_{+\mathrm{ext}}`$ is the smaller root of Eq. (17), and again we have only an intermediate stable phase.
In the grand-canonical ensemble, the stability analysis can be carried out by calculating the determinant of Hessian matrix of the energy with respect to $`S`$ and $`Q`$. The zeros of the determinant of Hessian matrix are given by $`St_{\mathrm{gc}}=0`$, where $`St_{\mathrm{gc}}`$ is
$`St_{\mathrm{gc}}`$ $`=`$ $`(d2)[(d7)\alpha ^3+2r_{+\mathrm{st}}^2(d10)\alpha ^218r_{+\mathrm{st}}^4\alpha +3r_{+\mathrm{st}}^6(d3)]`$ (31)
$`+18r_{+\mathrm{st}}^{2d+10}q^2\alpha 6r_{+\mathrm{st}}^{2d+12}q^2`$
Three cases happen for the roots of $`St_{\mathrm{gc}}=0`$:
1. It has two real roots $`r_{+\mathrm{st1}}<r_{+\mathrm{st2}}`$. In this case, the determinant of Hessian matrix is positive provided the radius of horizon lies in the range between $`r_{+\mathrm{larger}}<r_+<r_{+\mathrm{st2}}`$, where $`r_{+\mathrm{larger}}`$ is the larger value between $`r_{+\mathrm{ext}}`$ and $`r_{+\mathrm{st1}}`$. Thus, there exist an intermediate stable phase
2. It has one real root $`r_{+\mathrm{st2}}`$. In this case, there exists an intermediate stable phase with radius of horizon between $`r_{+\mathrm{ext}}`$ and $`r_{+\mathrm{st2}}`$.
3. It has no real solution. In this case the black hole is unstable for the whole range of $`r_+`$.
Numerical analysis shows that $`r_{+\mathrm{st2}}<r_{+\mathrm{st}}`$ (if $`r_{+\mathrm{st2}}`$ exists), and therefore the region of stability is smaller for the grand-canonical ensemble. This is due to the fact that the number of thermodynamic variables in the canonical ensemble is less than that of the grand-canonical ensemble.
## V General Solutions
Finally we give the general solutions of third order Lovelock gravity in $`d`$ dimensions for any arbitrary values of $`\alpha _2`$ and $`\alpha _3`$. In this case the solution of Eq. (13) is
$$f(r)=1+r^2\frac{(\frac{\alpha _2}{6\alpha _3})}{(_2^{d5})}+\frac{r^{\frac{2(d5)}{3}}}{\alpha _3}\left(\xi _d+\frac{\alpha _3}{24}\sqrt{\zeta _d}\right)^{1/3}+\frac{r^{\frac{2(d+1)}{3}}\left(\frac{1}{(_2^{d5})^2}\alpha _2^2\frac{1}{2(_4^{d3})}\alpha _3\right)}{36\alpha _3\left(\xi _d+\frac{\alpha _3}{24}\sqrt{\zeta _d}\right)^{1/3}}$$
(32)
where $`\xi _d`$ and $`\zeta _d`$ are
$$\xi _d=\frac{\alpha _2^3r^{2(d2)}}{216(_2^{d5})^3}\frac{\alpha _2\alpha _3r^{2(d2)}}{144(_4^{d3})(d6)(d5)}+\frac{m\alpha _3^2(d2)r^{d3}}{120(_5^{d2})}+\frac{\alpha _3^2q^2}{120(_5^{d2})(d3)}$$
$`\zeta _d`$ $`=`$ $`{\displaystyle \frac{\alpha _3r^{4(d2)}}{648(_4^{d3})^3}}{\displaystyle \frac{\alpha _2^2r^{4(d2)}}{36(d6)^2(d5)^2(_4^{d3})^2}}+{\displaystyle \frac{2^{\frac{1+(1)^d}{2}}\alpha _3^2mq^2(d2)r^{d3}}{25(d3)(_5^{d2})^2}}`$
$`{\displaystyle \frac{\alpha _2\alpha _3m(d2)r^{3d7}}{15\times 2^{1(1)^d}(d6)(d5)(_5^{d2})(_4^{d3})}}+{\displaystyle \frac{\alpha _3^2m^2(d2)^2r^{2(d3)}}{25\times 4^{1(1)^d}(_5^{d2})^2}}`$
$`{\displaystyle \frac{\alpha _2\alpha _3q^2r^{2(d2)}}{3(_4^{d3})(_4^{d2})(d3)(d5)(d6)^2}}+{\displaystyle \frac{8\times 2^{\frac{1+(1)^d}{2}}(d2)m\alpha _2^3r^{3d7}}{45(d5)^3(d6)^3(_5^{d2})}}`$
$`+{\displaystyle \frac{q^4\alpha _3^2}{25(d3)^2(_5^{d2})^2}}+{\displaystyle \frac{16q^2\alpha _2^3r^{2(d2)}}{9(d3)(d2)(d6)^3(d5)^3(_4^{d3})}}`$
In the above equations $`m`$ is an integration constant which is related to the mass of the solution. This solution is asymptotically flat and has a curvature singularity at $`r=0`$. Note that this asymptotically flat solution has three fundamental constants $`G=1`$, $`\alpha _2`$ and $`\alpha _3`$. Numerical analysis shows that the solutions (32) may present black holes with two inner and outer horizons, extreme black holes, black holes with one event horizon, or naked singularities depending on the values of $`\alpha _2`$, $`\alpha _3`$, $`m`$ and $`q`$. To be more clear we give the range of values of $`m`$, $`\alpha _2`$ and $`\alpha _3`$ for the 7-dimensional uncharged solution in order to have black hole solutions. For $`\alpha _3=3\alpha _2/2`$, any solution with $`m`$ and $`\alpha _2`$ in region I of Fig. 1 represents a black hole solution with an event horizon.
Unlike to the case of solutions with special values of $`\alpha _2`$ and $`\alpha _3`$ given in Sec. III, the uncharged solution for arbitrary values of fundamental constants $`\alpha _2`$ and $`\alpha _3`$ can have horizon even for $`d>9`$. Figure 2 shows the function $`f(r)`$ for $`d=10`$ as a function of $`r`$. It shows that there is an event horizon at $`r=2`$.
To investigate the first law of thermodynamics for the solutions with arbitrary values of $`\alpha _2`$ and $`\alpha _3`$, one should perform a numerical analysis. Using the expression for the entropy and the charge given in Eqs. (23) and (24), and the fact that $`M=[\mathrm{\Sigma }_{d2}/(16\pi )]m`$, where $`m`$ is the solution of $`f(m,r_+)=0`$, one can show numerically that the temperature and electric potential of Eq. (28) are exactly equal to the temperature and electric potential which can be computed through the use of geometrical analysis. Thus, the first law of thermodynamics is valid for the general solution given in Eq. (32). Also, one can perform a stability analysis in canonical and grand-canonical ensembles numerically. Again, as in the case of special solutions of third order Lovelock gravity, numerical analysis shows that there exist only an intermediate stable phase for black hole solutions which is smaller in grand-canonical ensemble .
## VI CLOSING REMARKS
In this paper, we added the third order Lovelock terms to the Gauss-Bonnet-Maxwell action, and introduced a new class of static solutions which are asymptotically flat, AdS or dS for $`\mathrm{\Lambda }=0`$, $`\mathrm{\Lambda }<0`$ or $`\mathrm{\Lambda }>0`$ respectively. We were only interested in the asymptotically flat solutions.
First, we consider the solutions for special values of $`\alpha _2`$ and $`\alpha _3`$ given in Eq. (14). For uncharged solutions, we found that for $`\alpha >0`$, the solutions present black holes provided the mass is larger than a critical value, $`m_{\mathrm{crit}}`$, given in Eq. (16). For negative $`\alpha `$, they present black holes with two inner and outer horizons for $`m>m_{\mathrm{ext}}`$, extreme black holes for $`m=m_{\mathrm{ext}}`$ or naked singularity for $`m<m_{\mathrm{ext}}`$, where $`m_{\mathrm{ext}}`$ is given in Eq. (18). These kind of uncharged solutions exist only in third order Lovelock gravity and do not happen in Einstein or Gauss-Bonnet gravity. For the case of charged solutions, we found that there exists an extremal value of $`r_+`$, given by Eq. (19), that determines whether the solution presents a black holes with two inner and outer horizons, an extreme black hole or a naked singularity. Accordingly we obtained temperature, entropy, charge , electric potential and mass of these black hole solutions. We also investigated the first law of thermodynamics, and found that these thermodynamics quantities satisfy the first law of thermodynamics. Also, we performed a stability analysis in canonical ensemble by considering the heat capacity of the solution and found that there exist two unstable phases separated by an intermediate stable phase. This analysis was also done through the use of the determinant of the Hessian matrix of $`M(S,Q)`$ with respect to its extensive variables and we got the same phase behavior with a smaller region of stability in the grand-canonical ensemble or no stable phase depending on the values of black hole’s parameters.
Second, we introduced the general solution of third order Lovelock gravity for arbitrary values of $`\alpha _2`$ and $`\alpha _3`$ and investigate its properties. Numerical analysis showed that these solutions may be interpreted as black hole solutions with two inner and outer event horizons, extreme black holes or naked singularity depending on the parameters of the solutions. We also found that the conserved and thermodynamics quantities for these general solutions satisfy first law of thermodynamics. Again, the stability analysis showed that there exists only an intermediate stable phase, for the black hole solutions which is smaller for the grand-canonical ensemble.
As we mention, the asymptotically flat solutions obtained in Chr contain only one fundamental constant. Finding new solutions in continued Lovelock gravity with more fundamental constants remains to be carried out in future. Also, the generalization of these solutions to the case of rotating solutions will be given elsewhere.
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# Budded membrane microdomains as regulators for cellular tension
## Abstract
We propose a mechanism for the control of the mechanical properties of the membrane of living cells that is based on the exchange of membrane area between the cell membrane and a membrane reservoir. The reservoir is composed of invaginated membrane microdomains which are liable to flatten upon increase of membrane strain, effectively controlling membrane tension. We show that the domain shape transition is first order, allowing for coexistence between flat and invaginated domains. During coexistence, the membrane tension is controlled by the domains elasticity and by the kinetics of the shape transition. We show that the tension of the plasma membrane of living cells is inherently transient and dynamical, and that valuable insights into the organization of the cell membrane can be obtained by studying the variation of the cell membrane tension upon mechanical perturbation.
Endocytosis, exocytosis, cell motility and many other crucial cellular processes are known to be influenced by the tension of the cell membranesheetz . The force needed to pull a tubular membrane tether from the cell with an optical trap is often used as a probe of the mechanical state of the membranesheetz\_cytoskel . The mechanical response of living cells to such perturbation involves cytoskeleton deformation, the breaking of membrane-cytoskeleton bondssheetz\_pip2 and changes in membrane morphology. The level of membrane tension in cells is thought to primarily reflects cytoskeleton anchoring to the membranesheetz\_cytoskel . Nevertheless, there is experimental evidence that Surface Area Regulation (SAR) occurs morris\_SAR and is able to buffer the increase of tension upon mechanical perturbation. SAR could be achieved by the transfer of membrane area from a reservoir to the plasma membrane, during which the level of membrane tension is at least partly controlled by the dynamical response of the reservoir. One possible manifestation of the exchange of membrane with a reservoir is the presence of a plateau in the force-extension curve of tether extraction experiments, as observed in sheetz\_buffer . The plateau ends when the putative membrane reservoir is emptied. In this particular experiment the reservoir is able to react quite fast to the perturbation ($`<0.1`$s), which could indicate i) that the reservoir is permanently connected to the plasma membrane and ii) that the regulation is purely physical and does not involve biological signaling.
Our goal in this paper is to study the influence of membrane morphology, and in particular the flattening of invaginated ”bud-like” domains, on the mechanical response of the membrane. Such domains include caveolae that form $`\mathrm{\Omega }`$ shaped invaginationscav\_review reactive to membrane stresscav\_stress ; woodman . Striking experiments, showing Caveolae flattening under tensionflat\_cav , support the idea that the delivery of invaginated membrane area to the plasma membrane is controlled by tensioncav\_reservoir . The shape of membrane domains in artificial membrane systems, such as giant vesicles, is also known to be dependent on the membrane tensionmembr\_domains1 . Indeed, domain budding has been observed upon the decrease of membrane stressmembr\_domains2 .
Our model is summarised in Fig.1. Here the cytoskeleton is included as linear elastic element (spring) mechanically coupled to the cell membrane. Although crude, this approximation is consistent with the experimental results of sheetz\_buffer (see SUPP ). Furthermore, it allows us isolate and study the response of the reservoir, which can be thought of as one element in a chain of response that transmits stress to the cytoskeleton. Cytoskeleton elasticity and membrane cytoskeleton anchoring are reflected by an effective stretching coefficient $`K_s`$ of order $`5.10^4k_BT/\mu m^4`$SUPP . Within this framework, the variation of the cell tension with the tether area $`𝒜_T`$ reads:
$$\gamma =\gamma _0+K_s\left(𝒜_T+𝒜_{res}𝒜_{res}^{(0)}\right)$$
(1)
where $`\gamma _0`$ and $`𝒜_{res}^{(0)}`$ are the cell tension and reservoir area at rest (without tether, $`𝒜_T=0)`$. If membrane area is delivered to the plasma membrane, as during exocytosis, the “tether” area is negative. One can see that the membrane tension $`\gamma `$ can be maintained constant upon tether pulling only if the decrease of reservoir area matches the increase of tether area.
The membrane domains constituting the reservoir are described within a very general theoretical framework. It is based on the domain bending rigidity, that disfavors the budded state, and its composition difference with the rest of the membrane, that promotes domain invagination in order to reduce the length of the domain periphery (Fig.2). Domains are treated as spherical caps of fixed area $`S`$ and adjustable curvature $`C`$. Their shape is uniquely characterized by the parameter $`\beta `$ (Fig.2 and Eq.(2)), equal to unity for a fully budded domain (a sphere), and that vanishes for a flat domain. The energy $`f`$ of a domain contains a surface tension-independent part $`f_{\sigma ,\kappa }`$. This itself involves a term arising from the line tension $`\sigma `$ membr\_domains1 , proportional to the length of the cap edge (neck), and a term giving the bending energy of the cap, proportional to the bending rigidity $`\kappa `$ and the squared curvaturesafran . Simple geometry then gives
$$f_{\sigma ,\kappa }[\beta ]=\sqrt{\pi S}\sigma \sqrt{1\beta }+8\pi \kappa \beta \beta =SC^2/16\pi $$
(2)
If the domain size exceeds a critical value $`S_c=\pi (4\kappa /\sigma )^2`$, an invaginated sphere ($`\beta =1`$) has a lower energy than a flat domain ($`\beta =0`$) and budding is expectedlipowsky\_budding . We assume that $`S>S_c`$ in what follows. Typical values of the parameters $`\sigma k_BT/nm`$ and $`\kappa =20k_BT`$ correspond to a critical size $`S_c=(100nm)^2`$, similar to the size of caveolae. In practice invaginated domains remain attached to the mother membrane by a small neck. Here, this is controlled phenomenologically by assigning the value $`\alpha _{bud}`$ ($`1`$) to the ratio of invagination radius $`R`$ to neck size, giving an upper bound $`\beta _{bud}=1(\alpha _{bud}/2)^2`$ to the shape parameter.
Including the membrane tension $`\gamma `$, the domain energy reads: $`f[\beta ]=f_{\sigma ,\kappa }[\beta ]+\gamma S\beta `$. Increasing membrane tension increases the energy of curved states $`\beta >0`$ and promotes the flat state (see Fig.2). The flat states eventually becomes stable for a critical tension $`\gamma ^{}`$:
$$S\gamma ^{}=(f_{\sigma ,\kappa }[1]f_{\sigma ,\kappa }[0])=2\sqrt{\pi S}\sigma 8\pi \kappa $$
(3)
As can be seen in Fig.2, the budded and flat domain shapes are separated by an energy barrier for intermediate tension. The existence of this barrier is crucial to their function as tension regulators, as it allows the coexistence of flat and invaginated domains. Variation of the membrane strain (the tether length) may occur without change of tension by adjusting the fraction of budded domains.
The energy scales in this system (with $`S=0.1\mu m^2`$) are
$$\overline{\sigma }\sqrt{\pi S}\sigma \overline{\kappa }8\pi \kappa 500k_BT;\overline{\gamma }\gamma S$$
(4)
The energy of surface tension competes with the line and bending energies for $`\overline{\gamma }500k_BT`$, or $`\gamma =\mathrm{2\; 10}^5J/m^2`$ (corresponding to tether forces of order $`10pN`$). This is precisely in the range of mechanical tension recorded for cellular membranessheetz\_cytoskel , which is very encouraging for the biological relevance of our model. One notes that the energy scale is very large compared to the thermal energy $`k_BT`$, or to the energy of any “active temperature” present in biological systemsmembrane\_activity . This has two important physical consequences: (i) the shape transition of a domain is very discontinuous, a domain snaps open rather than continuously flattening upon tension increase, and (ii) the budding and flattening transitions should actually occur at different tensions, for which the respective energy barriers are of order $`k_BT`$. In biological systems, the “temperature” $`T`$ might be seen as a parameter reflecting cellular activity, such as the polymerization of the actin cortex near the membrane, and the activity of membrane pumps. For simple cells such as Red Blood Cells, it is typically a few times the thermodynamic temperaturemembrane\_activity .
The bottleneck for the shape transition is the maximum of energy, which corresponds to a shape parameter $`\beta _{max}=1(\overline{\sigma }/(\overline{\kappa }+\overline{\gamma }))^2`$. The budding and flattening tensions ($`\gamma ^{(1)}`$ and $`\gamma ^{(0)}`$ respectively) at which the corresponding energy barrier vanishes are:
$$\overline{\gamma }^{(1)}=\overline{\sigma }\overline{\kappa }<\overline{\gamma }^{}=2\overline{\sigma }\overline{\kappa }<\overline{\gamma }^{(0)}=2\alpha _{bud}\overline{\sigma }\overline{\kappa }$$
(5)
where $`\alpha _{bud}`$ characterize the size of the invagination neck (see above).
We investigate the tension regulation performed by a collection of $`𝒩`$ domains, of total area $`𝒩S\beta _{bud}`$ (where $`\beta _{bud}=1(\alpha _{bud}/2)^2`$, is the largest value of shape parameter consistent with the existence of a finite size neck).
When flat and budded domains coexist, a fraction $`ϵ`$ of domains are invaginated, and the reservoir area is $`𝒜_{res}=S𝒩ϵ\beta _{bud}`$. The total membrane energy, including the contribution $`f_{\sigma ,\kappa }`$ of each of the $`𝒩`$ membrane domains (Eq.(2)) and the total work done against membrane tension can be written:
$$=𝒩\left(ϵf_{\sigma ,\kappa }[\beta _{bud}]+(1ϵ)f_{\sigma ,\kappa }[0]\right)+𝑑𝒜\gamma [𝒜]$$
(6)
Optimizing the energy for the fraction of invaginated domains $`/ϵ=0`$ leads directly to the regulation of membrane tension, when flat and budded domains coexist ($`0<ϵ<1`$). Substituting in Eq.(6) the expression for the surface tension Eq.(1) (with $`\beta _{bud}1`$), we find that the tension is set to the value $`\gamma ^{}`$ of Eq.(3), which depends on the characteristics of the membrane reservoir ($`\sigma `$, $`\kappa `$, and $`S`$), but not on the tether area. Regulation is achieved by adjusting the fraction of budded domains to:
$$ϵ^{}=ϵ_0\frac{(\overline{\gamma }_0+\overline{K_s}𝒜_T/S)\overline{\gamma }^{}}{\overline{K_s}𝒩}\overline{K_s}K_sS^2$$
(7)
where $`ϵ_0`$ is the fraction of budded domains corresponding to the tension at rest $`\gamma _0`$, and where a normalized stretching coefficient $`\overline{K_s}(0.1\overline{\sigma })`$, with dimension of energy, is introduced for convenience.
If the reservoir is given time to equilibrate, regulation starts for a level of perturbation corresponding to a tether area $`𝒜_T^{(1)}`$ (at which all domains are budded, $`ϵ^{}=1`$), and ends at a tether area $`𝒜_T^{(0)}`$ (at which all domains are flat, $`ϵ^{}=0`$). The tension of the cell membrane is then set to the values $`\gamma ^{}`$, for any perturbation within the range $`𝒜_T^{(1)}<𝒜_T<𝒜_T^{(0)}`$. If the perturbation is very fast, one expect a large difference between the regulated tension upon tether pulling and tether retraction, in agreement with Eq.(5).
To obtain the full kinetic response of the membrane to strain, we describe the transition as a classical Kramers’ processvankampen , were the transition time between two states is exponential with the energy barrier $`\mathrm{\Delta }f`$ that has to be overcome in the process: $`\tau =\tau _0\mathrm{exp}[\mathrm{\Delta }f/k_BT]`$. Here, $`\tau _0`$ is the characteristic fluctuation time of the domain shape, and $`T`$ may be an effective temperature resulting from the large biochemical activity near the cell membranemembrane\_activity . The transition time is very much dependent upon the membrane tension. Assuming that the transition of a single domain occurs with negligible change of tension (this implies $`𝒩1`$), the transition is fully described by the energy $`f[\beta ]=f_{\sigma ,\kappa }[\beta ]+\overline{\gamma }\beta `$, with $`\gamma `$ given by Eq.(1). Assuming for simplicity that both domain flattening and budding involve the same fluctuation time $`\tau _0`$, the kinetic evolution of the fraction $`ϵ`$ is given by
$$\tau _0\frac{dϵ}{dt}=ϵe^{\frac{f_{max}f[\beta _{bud}]}{k_BT}}+(1ϵ)e^{\frac{f_{max}f[0]}{k_BT}}$$
(8)
where the maximum of energy $`f_{max}`$ corresponds to the least favorable domain shape $`\beta _{max}=1\overline{\sigma }^2/(\overline{\kappa }+\overline{\gamma })^2`$
In order to mimic the tether pulling experiment, where the tether is typically extracted at constant speed ($`4\mu m/s`$) in sheetz\_buffer , we consider the reservoir response to a perturbation applied with a given rate $`\dot{𝒜}_T`$: $`𝒜_T=𝒜_T^{(0)}+\dot{𝒜}_Tt`$. If the perturbation is applied slowly ($`\dot{𝒜}_T\tau _0𝒜_{res}`$), the reservoir has time to equilibrate ($`dϵ/dt=0`$), and the fraction $`ϵ_T^{}`$ of budded domain is found from Eq.(8) to be given by $`ϵ_T^{}/(1ϵ_T^{})=e^{(f[0]f_{bud})/(k_BT)}`$, with$`f[0]f_{bud}2\overline{\sigma }(\overline{\gamma }(ϵ)+\overline{\kappa })`$. The fraction $`ϵ_T`$ is the equivalent of the equilibrium fraction $`ϵ^{}`$ (Eq.(7)), that takes thermal fluctuations into account ($`ϵ_T^{}ϵ^{}`$ if $`\overline{\sigma }k_BT`$). Thermal fluctuations smoothen the transition between budded and flat domains by allowing states of non-minimal energy to be populated. As a consequence, the tension is not perfectly constant during the transition, and the slope at mid plateau is of order $`\gamma ^{}/𝒜_T{}_{|_{plat}}{}^{}4k_BT/(S𝒜_{res})`$. If, on the other hand, the perturbation is applied very fast, the shape transition requires small energy barriers, which means high tension for bud flattening, and low tension for domain budding, respectively $`\gamma ^{(0)}`$ and $`\gamma ^{(1)}`$ give by Eq.(5).
The physical mechanism at the origin of tension regulation and the membrane hysteretic response to tether extraction and retraction are shown in Fig.3. To obtain an analytical expression of the plateau height with the perturbation rate, we approximate that the tension is almost constant during regulation ($`d\overline{\gamma }/dt0`$) so that the energy barrier is of order $`\mathrm{\Delta }f/(k_BT)\mathrm{log}𝒜_{res}/(\dot{𝒜}_T\tau _0)`$. The plateau tension upon increase and decrease of the perturbation are then respectively given by $`\overline{\gamma }_{10}=\overline{\gamma }^{(0)}2\alpha _{bud}^{}{}_{}{}^{3/2}\sqrt{\overline{\sigma }k_BT\mathrm{log}}`$ and $`\overline{\gamma }_{01}=\overline{\gamma }^{(1)}+\sqrt{\overline{\sigma }k_BT\mathrm{log}}`$, with $`\mathrm{log}\mathrm{log}[𝒜_{res}/(2\dot{𝒜}_T\tau _0)]`$. As expected for activated processes, the dependence of the tension at transition with the rate of perturbation $`\dot{𝒜}_T`$ is logarithmic. The same is true for the slope of the plateau, which can be estimated by identifying the plateau inflexion point. The condition $`d^2\overline{\gamma }/dt^2=0`$ imposes $`\dot{\mathrm{\Delta }f}/(k_BT)\dot{ϵ}/ϵ`$, corresponding to a plateau slope $`\overline{\gamma }_{10}/𝒜_T{}_{|_{plat}}{}^{}\sqrt{2k_BT\overline{\sigma }\alpha _{bud}^{}{}_{}{}^{3}/\mathrm{log}}/𝒜_{res}`$. As the effective membrane temperature may be inferred independently from measurement of membrane fluctuationsmembrane\_activity , the study of both the height and slope of the tension plateau gives valuable information of the kinetics of area transfer between the reservoir and the plasma membrane.
The difference between the quasi-static and dynamic plateau can be of order $`10^4J/m^2`$ and correspond to a difference in force of order $`10pN`$. This is precisely the scale of the forces measured upon tether extraction in sheetz\_buffer . This possibility thus exists that the initial increase of tension observed prior to the plateau in sheetz\_buffer is of purely kinetic origin and originate from the slow response of an already partially unfolded reservoir at rest. In this case, the membrane tension of the cell at rest would be $`\gamma ^{}`$ of Eq.(5), fully controlled by the mechanical properties of the membrane domains forming the reservoirSUPP .
In summary, we have derived the mechanical reactivity of a cell membrane, in contact with a reservoir composed of invaginated membrane domains liable to flatten under strain. The flattening transition is first order, which means that the invaginations snap open above a critical strain rather than continuously flattening, an observation consistent with experimental evidences from the membrane invagination caveolaecav\_reservoir . As a consequence, the cell mechanical response shows a plateau during the transition, corresponding to the coexistence of flat and invaginated domains. This study provides the basis for a mechanical regulation of the tension of the cell membrane. The mechanism at the origin of this regulation is of purely physical origin, consistent with the fast timescale ($`<0.1sec`$) of the observed cellular response. In practice, a ring of specialized membrane proteins such as dynamincav-dynamin is often present at the neck of membrane invagination. These proteins most probably influence the domain line energy, and might even dominate the energy required to flatten the domain. The regulation of membrane tension by membrane invaginations rely on the existence of two well defined domain shapes, separated by an energy barrier. If anything, neck proteins can only increase the energy barrier to flattening, thereby reinforcing tension regulation.
This work also opens the possibility of a new quantitative “force spectroscopy” of the cell membrane. One could thereby obtain structural information on the membrane organization, in much the same way information on a protein structure can be gathered from force measurement upon protein unfoldingprotein\_unfolding . As a first step, one may identify the fairly regular oscillation of the force during regulation in sheetz\_buffer , to the flattening of single domains. Preliminary analysisSUPP hints at domains of area $`S(400nm)^2`$.
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# References
>
Symbology from set theory applied to ecological systems: Gause’s exclusion principle and applications
J. C. Flores<sup>a,b</sup>
<sup>a</sup>Universidad de Tarapacá, Departamento de Física, Casilla 7-D, Arica-Chile
<sup>b</sup>Centro de Estudios del Hombre en el Desierto, CIHDE, Casilla 7-D, Arica-Chile
Abstract: We introduce a symbolic representation like set theory to consider ecologic interactions between species (ECOSET). The ecologic exclusion principle (Gause) is put in a symbolic way and used as operational tool to consider more complex cases like interaction with sterile species (SIT technique), two species with two superposed sources (niche differentiation) and N+P species competing by N resources, etc. Displacement (regional or characters) is also considered by using this basic tool. Our symbolic notation gives us an operative and easy way to consider elementary process in ecology. Some experimental data (laboratory or field) for ecologic process are re-considered under the optic of this set-theory.
Keywords: Coexistence and competition; Food web theory; Ecology; Set theory.
I Introduction.-
Interactions between species in ecology is after some time the object of study of mathematical branches (Gertsev et al (2004) and references therein). For instance, mathematical models for predator-prey are usually treated with nonlinear coupled differential equations like Lotka-Volterra (Begon et al (1999), Murray (1993), and references therein). It is true that, as a general rule, modellation of ecological systems is a difficult task since complexity in biological sciences is almost always present. For instance, the more known application of Lotka-Volterra system, that is the Hare-Lynx predator-prey data recorder by Hudson Bay Company in 1953, presents some troubles. The expected periodic solution between prey and predator is not determined from data (Murray (1993)). The reasons are not clear but one can expect a complex dynamics for real predator-prey system in a not isolated region. Very refined experiment in laboratories are actually realized (Costantino et al (1995)), nevertheless the more impressive laboratory experiments were carried-out some time ago by ecologist G. F. Gause with protozoan Paramesium species. These experiments have determined some general principles applied to inter-competition and coexistence in ecology. The more basic for us is the statement that: when we have two species competing by (exactly) the same niche then one of them disappear (Gause (1934)). That is, finally, one species full the niche. The statement is very restrictive since it requires quite special conditions (section IV) and in this paper it will be used as a basic principle to consider more complex situations.
Mathematical models like Lotka-Volterra for two species in competition predict coexistence for some range of parameters (high intraspecific competition). In fact, any mathematical system with atractors (out-side of the axis) in this two dimensional phase-space makes the same prediction. So, for us this mathematical result is not proof of real violation of the principle. Gause’s exclusion principle will be assumed for interspecific competition as a basic statement in this paper. Nevertheless, we accord that the distinction between the minimum amount of niche differentiation (in real ecological system) to produce its break is a difficult point to consider in practice. For an appropriate discussion see reference (Begon et al (1999)).
In this paper we will consider a mathematical modeling for ecological systems but using a symbology, called ECOSET, similar to basic set theory. This schema has the advantage of a condensed notation for a variety of ecological interacting systems. There are some similarities with usual set theory but also some differences. For instance, Gause’s exclusion principle (16) has not equivalence in the usual set theory. In our construction, Gause’s exclusion principle will be used as a basic operational tool to consider complex situations like more than one species and more than one resource. A great part of field and laboratory examples found in this article were hold from reference Begon et al (1999).
In this paper, species will be represented by capital letter like $`A`$, $`B,`$…etc. We will use the symbol $`S`$ only for a primary source at the basis (consume) of a ecological chain. We will give to $`S`$ a stronger sense: it will design a defined ensemble of sources which make viable the development of species. Namely, it is considered in the sense of an ecological niche. We assume that this primary resource is auto-sustained or externally sustained. The absence of species in a given region will be denoted by $`\varphi `$ in analogy with the symbol of usual set theory.
II Symbology for depredation ($`>`$) and basic definitions.-
Consider a primary resource $`S`$ and species $`A,B,C\mathrm{}\varphi `$. The notation
$$A>B\text{ (}B\text{ consumes }A\text{),}$$
(1)
means: species $`B`$ exploits species $`A`$ as a resource in a sense of depredation. In this way, a basic chain becomes for instance
$$S>A>B>\varphi (\text{a basic depredation chain }),$$
(2)
namely, species $`B`$ exploits $`A`$ as a resource and, $`A`$ exploits the primary resource $`S`$. Note that the species $`\varphi `$ at the end of the chain means that $`B`$ is not prey for others.
For two no-interacting species $`A`$ and $`B`$ (also for ecological process) in a given ecological region we write $`AB`$. When two species ($`A`$ and $`B`$) consume the same primary resource ($`S`$), in a not depending way, we write
$$S>\left(AB\right),\text{ (independent depredation).}$$
(3)
When two species ($`A`$ and $`B`$) exploit the same resource ($`S`$) in a interdependent way we write
$$S>\left(AB\right)\text{ (interdepending depredation).}$$
(4)
The notation $`S>AB`$ (without braces) means: species $`A`$ consumes $`S`$ with the help of species $`B`$.
Note that when one of the interdependent species if $`\varphi `$ then, after a time, no depredation on $`S`$ must occur. In a symbolic way,
$$\left\{S>\left(A\varphi \right)\right\}S>\varphi ,$$
(5)
where the symbol $``$ has a temporal interpretation or, it defines a temporal direction. Namely, if we have the ecological system $`X`$ then, after a time we have $`Y`$ ($`X`$ $`Y`$). For instance, the ecological proposition: when one species $`A`$ has not source for depredation then it dies, could be written as:
$$\left\{S>\varphi >A\right\}\left\{S>\varphi \right\}.$$
(6)
We also define the symbol $``$ which will be interpreted as equivalence between ecological process. For instance, in the process (3) we always assume implicitly the equivalence:
$$\left\{S>\left(AB\right)\right\}\left\{\left(S>A\right)\left(S>B\right)\right\}.$$
(7)
Two notes: (a) A ecological process like $`S>A>B`$ does not mean that $`S>B`$. If it is true that also $`B`$ consumes $`S`$ we must write $`\left(S>A>B\right)\left(S>B\right)`$. (b) We have the equivalent notation $`A>BB<A`$.
As a field example consider the food web with four trophic levels from New Zealand stream community (Begon et al (1999) page 836). This ecological systems is composed by Algae ($`A`$), Herbivorous insects ($`H`$), Predatory insects ($`P`$) and Brown trout ($`B`$). The web food is :
$$\left\{A>H>P>B\right\}\left\{A>H>B\right\},$$
(8)
and then
$$(8)A>\left\{\left(H>P\right)H\right\}>B.$$
(9)
So, $`A`$ is a primary resource and $`B`$ the final predator. Note that if $`P\varphi `$ (extinction), the web food does not disappear completely since $`\left\{A>H>B\right\}`$. As expected, biodiversity leads stability of ecological systems.
As another field example consider the ecological trophic level at the Lauca National Park (Arica-Chile). There is a biodiversity group represented by ancient flora and fauna under extreme climatic condition (3.0 to 4.5 Km of altitude). The highly adapted species conform almost a close system. Particularly we have a partial food-chain composed by different species like: basic herbs, including the so-called bofedal, ($`H`$); a kind of camel called Vicunas ($`VC`$); a kind of rodent called Vizcacha ($`VZ`$); two predators (Puma ($`P`$) and Zorro ($`Z`$)). Also we have the carrion-eat species condor ($`K`$). A basic process of Lauca National Park is represented by
$$H>\left(VCVZ\right)>\left\{\left(PZ\right)K(PZ)\right\}.$$
(10)
In fact, species $`P`$ and $`Z`$ are natural competitors in this region. Naturally, the web is more complex of that represented by (10), for instance species $`K`$ also depends on natural dead of $`VC`$ and $`VZ`$; but is only a basic notation example for us.
III Symbology for competition ($``$) and basic definitions.-
To consider species in competition (no depredation) we will use the symbol $``$ (see later). Here we give some basic definitions, for instance, consider species $`A`$ and $`B`$ in struggle for some source like water, space, etc. The symbol:
$$AB\text{ (}A\text{ is perturbed by }B\text{),}$$
(11)
means that species $`B`$ perturbs (interferes) $`A`$. Note that for depredation we use other symbol ($`>`$). The above process could also written as $`BA`$ or in the equivalence language
$$AB\text{ }BA\text{.}$$
(12)
The symbols $`,,`$ and $``$ are used in the same way that in the above section.
With this basic definitions we can represent competition between species. In fact, if $`A`$ and $`B`$ are two species in competition (no depredation) we write
$$\left(AB\right)\left(BA\right),\text{ (}A\text{ and }B\text{ compete).}$$
(13)
By simplicity we will use the alternative symbol $``$ for competition. Namely,
$$\{(AB)(BA)\}\{AB\},\text{ (symbol for competition).}$$
(14)
Before to ending this section we give a useful definition which will be used in some cases. The notion of “potential competitors” is related to two species who put together then compete. In our symbolic notation, it could be written as
$$\{S>(AB)\}\{S>(AB)\}\text{ (potential competitors).}$$
(15)
IV Gause’s exclusion principle for interspecific competition and symbolic notation.-
Gause’s exclusion principle (or competitive exclusion principle) in ecology states that: when we have two species $`A`$ and $`B`$ which compete (interspecific competition) for the same invariable ecological primary resource $`S`$ (realized niche), then one of them disappear (Begon et al (1999), Gause (1934), Hasting (1996), Flores (1998)). It is important to note that Gause’s exclusion principle holds when no migration, no mutation and no resource differentiation exist in the ecological systems. Note that it refers to interspecific competition. The case of intraspecific competition will be touched briefly in section IX. The principle assures that the more stronger species in the exploitation of primary resource survives. Applications could be found in many text of ecology. A direct application of this principle to Neanderthal extinction in Europe could be found in reference (Flores (1998)).
In our symbolic notation the principle could be written as:
$$\{S>(AB)\}\text{ }\left\{(S>A)\text{ or }(S>B)\right\},\text{ (Gause).}$$
(16)
The above statement (16) will be a basic operational tool to consider more general cases or application like two sources and two predators, or more general. So, (16) is our start-point. The logic operator $`or`$ (some times written as $``$) is the usual exclusion symbol in set theory.
The more famous example of exclusion comes from the classic laboratory work of ecologist G. F. Gause (1934), who considers two type of Paramecium, namely, P. caudatum and P. aurelia. Both species grow well alone and reaching stable carrying capacities in tubes of liquid medium and consuming bacteria. When both species grow together, P. caudatum declines to the point of extinction and leaving P. aurelia in the niche.
As said before, other examples could be found in literature. For instance, competition between Tribolium confusum and Tribolium castaneum where one species is always eliminated when put together (Park (1954)).
V Application: interaction with sterile individuals and eradication (SIT).-
A corollary of the above principle can be found when one of the species in competition is sterile. In fact, we define a sterile specie $`M`$ as a species which exploits a resource $`S`$ and then disappear. Namely,
$$\left\{S>M\right\}\left\{S>\varphi \right\},\text{ (sterile species).}$$
(17)
Now, we consider this definition together to the exclusion principle. Let $`A`$ be a species which exploits the resource $`S`$, and let $`M`$ be a sterile species introduced which exploits the same resource. The application direct of the principle (16), and definition (17), tell us that
$$S>(AM)(S>A)\text{ or }(S>M),\text{ (Gause applied).}$$
(18)
and then,
$$\left\{S>A\right\}\text{ or }\left\{S>\varphi \right\},\text{ (}M\text{ is sterile ). }$$
(19)
Putting together (18) and (19), we have the ecological process:
$$\{S>(AM)\}\left\{(S>A)\text{ or }(S>\varphi )\right\},\text{ (Gause for sterile),}$$
(20)
so, at least one species of both disappear and then there is the possibility of total extinction in the niche ($`S>\varphi `$). The known SIT (Sterile Insect Technique, Barclay (2001)) uses this principle to eradicate undesirable insects. In fact, sterile insects compete with native ones for a source and there is the possibility of total extinction (eradication, $`S>\varphi `$). The fruit flies (medfly) eradication program carried out in many regions of the world, for instance in Arica-Chile, could be understand partially with the above results (Flores (2000) and (2003)). If $`S`$ is the female-native group then the native male group $`A`$ and the sterile male group $`M`$ compete by the “resource $`S`$”. In this way using (20) there is the possibility of $`S>\varphi `$ corresponding in this case to extinction of all type of male and then the wild species disappear.
VI Application: two species, two resources, and niche differentiation.-
As other application of our symbology for the principle of Gause, consider two resources $`S_1`$ and $`S_2`$ and two species $`A`$ and $`B`$ in competition by these resources. Namely, consider the ecological systems where
$$(S_1S_2)>(AB),$$
(21)
or
$$(21)\{S_1>(AB)\}\{S_2>(AB)\},$$
(22)
the equivalence becomes since both species consume any of two resources. From the exclusion principle (16), we have
$$(22)\left\{S_1>A\text{ or }S_1>B\right\}\left\{S_2>A\text{ or }S_2>B\right\},\text{ (Gause applied),}$$
(23)
and the four final possibilities:
(a) $`\left\{S_1S_2\right\}>A.`$ Species $`A`$ exterminates $`B`$.
(b) $`\left\{S_1S_2\right\}>B.`$ Species $`B`$ exterminates $`A`$.
(c) $`\left\{S_1>A\right\}\left\{S_2>B\right\}`$. Species $`A`$ exploits $`S_1`$ and $`B`$ exploits $`S_2`$.
(d)$`\left\{S_1>B\right\}\left\{S_2>A\right\}`$. Species $`A`$ exploits $`S_2`$ and $`B`$ exploits $`S_1`$.
The last two possibilities (c) and (d) tell us that both species could survive by resources exploitations in a differential (partitioned) way. Some time this coexistence is considered a violation of Gause’s exclusion principle; but it is not. In fact, we have more than one resource (realized niche).
The behavior found in the above process (a-d), has been observed in laboratories (see Begon et al (1999) page 311, or Tilman (1977)) where two diatom species (Asterionella formosa and Cyclotella meneghimiana) compete by silicate ($`S_1`$) and phosphate ($`S_2`$) as elementary resources. In fact, for different proportions of this components one can see extermination or stable coexistence (Tilman (1977)). This is a valuable laboratory experimental example which support our theory as a clear and efficient operational tool.
VII General case with $`N+P`$ species competing by $`N`$ resources.-
Considering the above result of section VI for two sources and two species in competition, it seems natural to extend it to a more general case. This will be do in this section. Consider $`N`$ primary resources $`S_i`$ ($`i=1,2,\mathrm{}N`$) and $`N+P`$ ($`P0`$) species $`A_j`$ ($`j=1,2,\mathrm{}N+P`$) competing by the resources. We will consider this species as potential competitors in the sense defined by expression (15). In this section, the mean result is that at least $`P`$ species disappear. So, we are considering the ecological systems given by the process
$$\left\{\underset{i=1}{\overset{N}{}}S_i\right\}>\left\{\underset{jk}{\overset{N+P}{}}(A_jA_k)\right\},$$
(24)
where the summation is understood in the sense of independent species in a region, explicitly, $`S_i=S_1S_2S_{3\mathrm{}.}`$ , the ecological process (24) is equivalent to
$$(24)\underset{i=1}{\overset{N}{}}\{S_i>\underset{jk}{\overset{N+P}{}}(A_jA_k)\},$$
(25)
because any species consumes any resources. Using Gause (16) for every pair $`i,k`$ then we have
$$(25)\underset{i=1}{\overset{N}{}}\left\{S_i>\left(A_1\text{ or }A_2\text{ or }A_3\text{…or }A_{N+P}\right)\right\}\text{, (Gause applied)}.$$
(26)
So every $`S_i`$ is consumed by one species; but note that one species could consume more than one resource. In this way, at least there are $`P`$ species extinct. The extrema option are:
(a) One species, called the exterminator, finally uses the $`N`$ resources.
(b) $`N`$ species coexist. That is, one species for every source (niche differentiation).
VIII Regional and character displacement.-
As said before, Gause competitive exclusion principle is a basic tool which could be applied to more complex cases as two resources or more. In this section we want to show how our notation is so coherent that it could be applied to other cases. In fact we will consider displacement of species. We will see it in a very operative way.
We define displacement of a species $`D`$ from a source $`S_1`$ to $`S_2`$ as
$$\left\{\left(S_1>D\right)\left(S_2>\varphi \right)\right\}\left\{\left(S_1>\varphi \right)\left(S_2>D\right)\right\},\text{ ( displacement).}$$
(27)
Note that displacement could be understood in two ways:
(a) Regional or spatial displacements (migration). Namely, $`S_1`$ and $`S_2`$ are sources in different spatial locations.
(b) Character displacements. Namely, $`S_1`$ and $`S_2`$ are sources in the same spatial place but species $`D`$ changes (displaces) its sources necessities, for instance due to mutation.
Consider two specie, $`A`$ and $`D`$, competing by the same resources $`S_1.`$ Assume that $`D`$ displaces to the unoccupied resource $`S_2`$ before to apply Gause. The ecological system is given by
$$\{S_1>(AD)\}\{S_2>\varphi \},$$
(28)
$$\left\{S_1>\left(DAAD\right)\right\}\left\{S_2>\varphi \right\},$$
(29)
$$\left\{S_1>DA\right\}\left\{S_1>AD\right\}\left\{S_2>\varphi \right\},$$
(30)
$$\left\{\left\{\left(S_1>D\right)\left(S_2>\varphi \right)\right\}A\right\}\left\{S_1>AD\right\},$$
(31)
where we have used $`\varphi A\varphi `$. In this stage, assuming that species $`D`$ displaces to $`S_2`$ (see (27)) we have
$$\left\{\left\{S_1>\varphi S_2>D\right\}A\right\}\left\{S_1>AD\right\},$$
(32)
$$\left\{S_1>\varphi A\right\}\left\{S_2>DA\right\}\left\{S_1>AD\right\},$$
(33)
using newly $`\varphi A\varphi `$ we obtain
$$\left\{S_2>DA\right\}\left\{S_1>AD\right\}.$$
(34)
In resume, from (28) and (34) we have:
$$\{S_1>(AD)\}\{S_2>\varphi \}\{S_1>AD\}\{S_2>DA\},$$
(35)
and both species survive due to displacement. A practical example for displacement comes from the same Gause classic experiments. In fact, when two protozoan P. caudatum and P. bursaria were grown together neither species suffered a decline to the point of extinction. They were in competition with one another but although they lived together in the same tube, they were spatially separated. P. caudatum lived suspended in the liquid medium and P. bursaria was concentrate at the bottom of the tube (Begon (1999). So, coexistence is related here with displacement.
An example where character displacement gives coexistence is provide by mud snails Hydrobia ulvae and Hydrolia ventrosa (Saloniemi (1993)). When they live apart, their sizes are almost identical. Nevertheless, when put together they reach different sizes in time. In fact, when they are similarly sized (apart) they consume similarly sized food; but when they are put together, the more larger tends to consume larger food particles.
IX Intraspecific competition : Gause does not hold.-
As mentioned in the introduction, an important and debated question is related to the validity of Gause’s principle when high intraspecific competition exist (individuals of the same species compete themselves by resources). So, we have two species $`A`$ and $`I`$ exploiting a resource $`S`$; but $`I`$ presents high intraspecific competition. In this case it seems that both species could coexist (Begon (1999)). That is, assume species $`A`$ is a weak consumer of $`S`$ and then in principle it must disappear face to the strong consumer $`I`$; but $`I`$ presents a so high degree of intraspecific competition that $`A`$ has a chance to survive.
Our formalism does not respond the question about coexistence, or not, in this case since Gause does not hold here. In fact, we define a species $`I`$ with intraspecific competition as
$$\left(S>I\right)\left(S>(II)\right),\text{ (intraspecific competition).}$$
(36)
Now, consider species $`I`$ competing with $`A`$ by the resource $`S`$, namely,
$$\left(S>IA\right)\left(S>AI\right),$$
(37)
since $`I`$ is (high) intraspecific competitor
$$(37)\left(S>\left(II\right)A\right)\left(S>AI\right),$$
(38)
and Gause does not hold since $`\left(II\right)I`$ (with exception of $`\varphi `$). Note that we use the term high intraspecific competition. This is so because in the above process we use intracompetition before applying Gause to process (37). Weak intraspecific means that in (37) we use Gause and after the intraspecific character of species $`I`$. In this case we have exclusion.
X Virtual process generation.-
In this section we consider a virtual (speculative) possibility to generate new processes from some known. In this sense, the new processes constructed are not necessarily real process. The advantage of this generation processes is to explore some future symmetries of ecological systems.
We define the dual ecological process of a given process as this one where the changes $`(>)()`$ and $`()(>)`$ operate. We define the inverse ecological process as this one where the changes $`(>)(<)`$ and $`()()`$ operate. For instance, the dual of $`\left(A>B\right)`$ denotes by $`\left(A>B\right)^D`$ is $`\left(AB\right),`$ namely, $`\left(A>B\right)^D\left(AB\right).`$ The inverse of $`\left(A>B\right)`$ is $`\left(A<B\right)`$, namely $`\left(A>B\right)^I\left(A<B\right)`$.
For instance, consider the ecological process of two species $`A`$ and $`B`$ in mutual depredation $`(><)`$ which consume also a primary resource $`S`$. Moreover, since individual of every specie dies, the primary resources uses this as a food resource (nutrients). So, consider the idealized process where
$$S>(A><B)>S,$$
(39)
obviously, this process is invariant under inversion operation. Namely,
$$\{S>(A><B)>S\}^I\{S>(A><B)>S\}.$$
(40)
As other hypothetical example, consider the process $`(S>BS)`$. From the above definition for dual and inverse operations we have $`(S>BS)^{DI}(S>BS)`$, namely, we have an invariant ecological process under dual and inversion operations.
XI Time for exclusion (number of generations).-
In this last section we will be concerned with a basic discussion of time for exclusion. That is, Gause is a time evolutive process; but it does not refer explicitly to how long is this time for exclusion. It seems quite natural to think that when more similar species in competition are, then more longer the extinction time becomes. To be more explicit, Consider two specie $`A`$ and $`B`$ in competition according with Gause. Assume species $`B`$ is excluded in a number $`N_B`$ of generations. We will assume that all similitude between both species could be quantified by one parameter $`s`$. That is, $`s=1`$ both species are completely similar (same species). Opposite, $`s=0`$, means that they are completely different species (genotype, phenotype, etc.). The two parameter $`N_B`$ (the extinction generation number) and the similitude parameter $`s,`$ are related by the simple expression $`N_B=1/(1s),`$ proposed originally in reference Flores (1998). So, more similar the species in competition are ($`s1`$), a much longer time (number of generation) is necessary to exclusion.
XI Conclusions.-
We have presented a symbology like to set theory applied to ecological interacting process (ECOSET). Chains of depredation or competition were explicitly studied. Particularly, Gause’s exclusion principle was considered in this notation and used as a basic operational tool. For instance, it was applied to competition with sterile individuals (SIT), two species with two resources and, more general, to $`N+P`$ species with $`N`$ resources. The symbology is so coherent that: displacement breaks exclusion was obtained with basic operations of our theory. Examples from laboratory and field were explicitly considered.
Resume for symbols
$``$ Perturbation (no depredation).
$`>`$ Depredation.
$``$ Two independent species in a region (eventually independent process).
$``$ Two interdependent species.
$``$ Temporal evolution.
$``$ Equivalence.
$``$ Abbreviation for competition.
$`><`$ Mutual depredation.
$`or`$ Exclusion (some times $``$).
$``$ Independent species (eventually process): $`ABCD\mathrm{}`$.
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# Rigidity of amalgamated product in negative curvature
## 1. Introduction
In , Y. Shalom proved the following theorem which says that for every lattice $`\mathrm{\Gamma }`$ in the hyperbolic space and for any decomposition of $`\mathrm{\Gamma }`$ as an amalgamated product $`\mathrm{\Gamma }=A_CB`$, the group $`C`$ has to be “big”. In order to measure how “big” $`C`$ is, let us define the critical exponent of a discrete group $`C`$ acting on a Cartan Hadamard manifold by
$$\delta \left(C\right)=inf\{s>0|\mathrm{\Sigma }_{\gamma \mathrm{\Gamma }}e^{sd(\gamma x,x)}<+\mathrm{}\}.$$
###### Theorem 1.1.
Let $`\mathrm{\Gamma }`$ be a lattice in $`PO(n,1)`$. Assume that $`\mathrm{\Gamma }`$ is an amalgamated product of its subgroups $`A`$ and $`B`$ over $`C`$. Then, the critical exponent $`\delta (C)`$ of $`C`$ satisfy $`\delta (C)n2`$.
An example is given by any $`n`$-dimensional hyperbolic manifold $`X`$ which contains a compact separating connected totally geodesic hypersurface $`Y`$. The Van Kampen theorem then says that the fundamental group $`\mathrm{\Gamma }`$ of $`X`$ is isomorphic to the free product of the fundamental groups of the two halves of $`XY`$ amalgamated over the fundamental group $`C`$ of the incompressible hypersurface $`Y`$. Such examples do exist in dimension $`3`$ thanks to the W.Thurston’s hyperbolization theorem. In any dimension, A. Lubotsky showed that any standard arithmetic lattice of $`PO(n,1)`$ has a finite cover whose fundamental group is an amalgamated product, cf. . In fact, A. Lubotsky proved that any standard arithmetic lattice $`\mathrm{\Gamma }`$ has a finite index subgroup $`\mathrm{\Gamma }_0`$ which is mapped onto a nonabelian free group. A nonabelian free group can be written in infinitely many ways as an amalgamated product, so one get infinitely many decomposition of $`\mathrm{\Gamma }_0`$ as an amalgamated product by pulling back the amalgamated decomposition of the nonabelian free group.
In these cases there is equality in theorem 1.1, ie. $`\delta \left(C\right)=n2`$ where $`C`$ is the fundamental group of $`Y`$, and Y. Shalom suggested in that the equality case in the theorem 1.1 happens only in that case.
The aim of this paper is to show that the theorem 1.1 still holds when $`\mathrm{\Gamma }`$ is the fundamental group of a compact riemannian manifold of variable sectional curvature less than or equal to $`1`$, and characterize the equality case.
###### Theorem 1.2.
Let $`X`$ be an $`n`$-dimensional compact riemanniann manifold of sectional curvature $`K1`$. We assume that the fundamental group $`\mathrm{\Gamma }`$ of $`X`$ is an amalgamated product of its subgroups $`A`$ and $`B`$ over $`C`$ and that neither $`A`$ nor $`B`$ equals $`\mathrm{\Gamma }`$. Then, the critical exponent $`\delta (C)`$ of $`C`$ satisfy $`\delta (C)n2`$. Equality $`\delta (C)=n2`$ happens if and only if $`C`$ cocompactly preserves a totally geodesic isometrically embedded copy $`^{n1}`$ of the hyperbolic space of dimension $`n1`$. Moreover, in the equality case, the hypersurface $`Y^{n1}:=^{n1}/C`$ is embedded in $`X`$ and separates $`X`$ in two connected components whose fundamental groups are respectively $`A`$ and $`B`$.
###### Remark 1.3.
(i) By the assumption on $`A`$ or $`B`$ not being equal to $`\mathrm{\Gamma }`$ we exclude the trivial decomposition $`\mathrm{\Gamma }=\mathrm{\Gamma }_CC`$ where $`A=\mathrm{\Gamma }`$ and $`C=B`$ can be an arbitrary subgroup of $`\mathrm{\Gamma }`$, for example any cyclic subgroup, in which case the conclusion of theorem 1.2 fails. Also note that because of this assumption on $`A`$ and $`B`$, we have $`AC`$ and $`BC`$ .
(ii) Let us recall that standard arithmetic lattices in $`PO(n,1)`$ have finite index subgroup with infinitely many non equivalent decompositions as amalgamated products, cf. . In fact, among these decompositions, all but finitely many of them are such that $`\delta (C)>n2`$. Indeed, by theorem 1.2, if $`\delta (C)=n2`$ then $`C`$ is the fundamental group of an embedded totally geodesic hypersurface in $`X`$, but there are only finitely many totally geodesic hypersurfaces by .
When a group is an amalgamated product, it acts on a simplicial tree without fixed point and theorem 1.1 is a particular case of the
###### Theorem 1.4.
(, theorem 1.6). Let $`\mathrm{\Gamma }SO(n,1)`$ , $`n3`$, be a lattice Suppose $`\mathrm{\Gamma }`$ acts on a simplicial tree $`T`$ without fixed vertex. Then there is an edge of $`T`$ whose stabilizer $`C`$ satisfies $`\delta (C)n2`$.
In the case $`\mathrm{\Gamma }`$ is cocompact, the conclusion of theorem 1.4 holds for the stabilizer of any edge which separates the tree $`T`$ in two unbounded components, and the proof of this is exactly the same as the proof of theorem 1.2. In particular, when the action of $`\mathrm{\Gamma }`$ on $`T`$ is minimal, (ie. there is no proper subtree of $`T`$ invariant by $`\mathrm{\Gamma }`$), the conclusion of theorem 1.4 holds for every edge of $`T`$, in the variable curvature setting, and we are able to handle the equality case.
###### Theorem 1.5.
Let $`\mathrm{\Gamma }`$ be the fundamental group of an $`n`$-dimensional compact riemannian manifold $`X`$ of sectional curvature less than or equal to $`1`$. Suppose $`\mathrm{\Gamma }`$ acts minimally on a simplicial tree $`T`$ without fixed point. Then, the stabilizer $`C`$ of every edge of $`T`$ satisfies $`\delta (C)n2`$. The equality $`\delta (C)=n2`$ happens if and only if there exist a compact totally geodesic hypersurface $`YX`$ with fundamental group $`\pi _1(Y)=C`$. Moreover, in that case, $`Y`$ with its induced metric has constant sectional curvature $`1`$.
Another interesting case contained in theorem 1.5 is the case of $`HNN`$ extension. Let us recall the definition of an $`HNN`$ extension. Let $`A`$ and $`C`$ be groups and $`f_1:CA`$, $`f_2:CA`$ two injective morphisms of $`C`$ into $`A`$. The $`HNN`$ extension $`A_C`$ is the group generated by $`A`$ and an element $`t`$ with the relations $`tf_1\left(\gamma \right)t^1=f_2\left(\gamma \right)`$. For example, let $`X`$ be a compact manifold containing a non separating compact incompressible hypersurface $`YX`$. Let $`A`$ be the fundamental group of the manifold with boundary $`XY`$ obtained by cuting $`X`$ along $`Y`$ and let $`C`$ be the fundamental group of $`Y`$. The boundary of $`XY`$ consists in two connected components $`Y_1XY`$ and $`Y_2XY`$ homeomorphic to $`Y`$. By the incompressibility assumption, these inclusions give rise to two embeddings of $`C`$ into $`A`$, and the fundamental group of $`X`$ is the associated $`HNN`$ extension $`A_C`$.
###### Theorem 1.6.
Let $`\mathrm{\Gamma }`$ be the fundamental group of a compact riemannian manifold $`X`$ of dimension $`n`$ and sectional curvature less than or equal to $`1`$. Suppose that $`\mathrm{\Gamma }=A_C`$ where $`A`$ is a proper subgroup of $`\mathrm{\Gamma }`$. Then, we have $`\delta (C)n2`$ and equality $`\delta (C)=n2`$ if and only if there exist a non separating compact totally geodesic hypersurface $`YX`$ with fundamental group $`\pi _1(Y)=C`$. Moreover, in that case, $`Y`$ with its induced metric is of constant sectional curvature $`1`$, and the $`HNN`$ decomposition arising from $`Y`$ is the one we started with.
Let us summarize the ideas of the proof of theorem 1.2. We work on $`\stackrel{~}{X}/C`$. The amalgamation assumption provides an essential hypersurface $`Z`$ in $`\stackrel{~}{X}/C`$, namely $`Z`$ is homologically non trivial in $`\stackrel{~}{X}/C`$. The volume of all hypersurfaces homologous to $`Z`$ is bounded below by a positive constant because their systole are bounded away from zero. We then construct a smooth map $`F:\stackrel{~}{X}/C\stackrel{~}{X}/C`$, homotopic to the identity which contracts the volume of all compact hypersurfaces $`Y`$ by the factor $`\left(\frac{\delta \left(C\right)}{n2}\right)^{n1}`$, namely $`vol_{n1}F\left(Y\right)\left(\frac{\delta \left(C\right)}{n2}\right)^{n1}vol_{n1}Y`$. This contracting property together with the lower bound of the volume of hypersurfaces in the homology class of $`Z`$ gives the inequality $`\delta \left(C\right)n2`$. This map is different from the map constructed in , in particular it can be defined under the single condition that the limit set of $`C`$ is not reduced to one point. Moreover, its derivative has an upper bound depending only on the critical exponent of $`C`$.
The equality case goes as follows. When $`\delta \left(C\right)=n2`$, the map $`F:\stackrel{~}{X}/C\stackrel{~}{X}/C`$ contracts the $`\left(n1\right)`$-dimensional volumes, ie. $`\left|Jac_{n1}F\right|1`$. This contracting property is infinitesimally rigid in the following sense. Let us consider a lift $`\stackrel{~}{F}`$ of $`F`$. If $`\left|Jac_{n1}\stackrel{~}{F}\left(x\right)\right|=1`$ at some point $`x\stackrel{~}{X}`$, then $`\stackrel{~}{F}\left(x\right)=x`$, there exists a tangent hyperplane $`ET_x\stackrel{~}{X}`$ such that $`D\stackrel{~}{F}\left(x\right)`$ is the orthogonal projector of $`T_x\stackrel{~}{X}`$ onto $`E`$ and the limit set $`\mathrm{\Lambda }_C`$ is contained in the topological equator $`E\left(\mathrm{}\right)\stackrel{~}{X}`$ associated to $`E`$. By topological equator $`E\left(\mathrm{}\right)\stackrel{~}{X}`$ associated to $`E`$, we mean the set of end points of those geodesic rays starting at $`x`$ tangently to $`E`$.
We then prove the existence of a point $`x\stackrel{~}{X}`$ such that
(1.1)
$$\left|Jac_{n1}\stackrel{~}{F}\left(x\right)\right|=1.$$
If there would exist a minimizing cycle in the homology class of $`Z`$ in $`\stackrel{~}{X}/C`$, any point of such a cycle would satisfy (1.1). As no such minimizing cycle a priori exists because of non compactness of $`\stackrel{~}{X}/C`$, we prove instead the existence of a $`L^2`$ harmonic $`\left(n1\right)`$-form dual to $`Z`$, which is enough to prove existence of a point $`x`$ such that (1.1) holds.
At this stage of the proof, there is a big difference between the constant curvature case and the variable curvature case.
In the constant curvature case, any topological equator bounds a totally geodesic hyperbolic hypersurface $`^{n1}`$, and therefore, as the group $`C`$ preserves $`\mathrm{\Lambda }_CE\left(\mathrm{}\right)=^{n1}`$, it is not hard to see that $`C`$ also preserves $`^{n1}`$ and acts cocompactly on it, and the hypersurface of the equality case in theorem 1.2 is $`^{n1}/C`$, .
In the variable curvature case, we first show the existence of a $`C`$-invariant totally geodesic hypersurface $`\stackrel{~}{Z}_{\mathrm{}}\stackrel{~}{X}`$ whose boundary at infinity coincides with $`\mathrm{\Lambda }\left(C\right)`$, and then we show that $`\stackrel{~}{Z}_{\mathrm{}}`$ is isometric to the real hyperbolic space. We then show that $`Y=:_{}^{n1}/C`$, which is compact, injects in $`X=\stackrel{~}{X}/\mathrm{\Gamma }`$ and separates $`X`$ in two connected components whose fundamental groups are $`A`$ and $`B`$ respectively.
In order to show the existence of such a totally geodesic hypersurface $`\stackrel{~}{Z}_{\mathrm{}}`$, we first prove that $`C`$ is a convex cocompact group, ie. the convex hull of the limit set of $`C`$ in $`\stackrel{~}{X}`$ has a compact quotient under the acton of $`C`$, and that the limit set of $`C`$ is homeomorphic to an $`\left(n2\right)`$-dimensional sphere.
The convex cocompactness property of $`C`$ and the fact that the limit set $`\mathrm{\Lambda }\left(C\right)`$ of $`C`$ is homeomorphic to an $`\left(n2\right)`$-dimensional topological sphere are the two key points in the equality case.
This compactness property then allows us to prove the existence of a mimimizing current in the homology class of the essential hypersurface $`Z\stackrel{~}{X}/C`$. By regularity theorem this minimizing current $`\stackrel{~}{Z}_{\mathrm{}}`$ is a smooth manifold except at a singular set of codimension at least 8. By the contracting properties of our map $`F`$, $`\stackrel{~}{Z}_{\mathrm{}}`$ is fixed by $`F`$ and the geometric properties of $`F`$ at fixed points where the $`\left(n1\right)`$-jacobian of $`F`$ equals 1 allows us to prove that $`\stackrel{~}{Z}_{\mathrm{}}`$ is totally geodesic and isometric to the hyperbolic space.
Let us now briefly describe the proof of the convex cocompactness property of $`C`$ in the equality case.
The group $`C`$ (or a finite index subgroup of it) actually globally preserves a smooth cocompact hypersurface $`\stackrel{~}{Z}\stackrel{~}{X}`$ which separates $`\stackrel{~}{X}`$ into two connected components and whose boundary $`\stackrel{~}{Z}\stackrel{~}{X}`$ coincides with $`\mathrm{\Lambda }_CE\left(\mathrm{}\right)`$. In the case where $`C`$ wouldn’t be convex cocompact, we are able to find an horoball $`HB\left(\theta _0\right)`$ centered at some point $`\theta _0\mathrm{\Lambda }_C`$ in the complementary of which lies the hypersurface $`\stackrel{~}{Z}`$.
The contradiction then comes from the following.
Consider a sequence of points $`\theta _i\stackrel{~}{X}`$ converging to $`\theta _0`$ and geodesic rays $`\alpha _i`$ starting from the point $`x\stackrel{~}{X}`$ at which $`\left|Jac_{n1}\left(x\right)\right|=1`$ and ending up at $`\theta _i`$. These geodesic rays have to cross $`\stackrel{~}{Z}`$ at points $`z_i`$ which are at bounded distance from the orbit $`Cx`$ of $`x`$, therefore the shadows $`𝒪_i`$ of balls centered at these $`z_i`$ enlighted from $`x`$ have to contain points of $`\mathrm{\Lambda }_C`$ by the shadow lemma of D. Sullivan. On the other hand, we show that it is possible to choose the sequence $`\theta _i`$ in such a way that these shadows $`𝒪_i`$ don’t meet $`\mathrm{\Lambda }_C`$. This property $`𝒪_i\mathrm{\Lambda }_C=\mathrm{}`$ comes from a choice of $`\theta _i`$ such that the distance between $`z_i`$ and the set $`H`$ of all geodesics rays at $`x`$ tangent to $`ET_x\stackrel{~}{X}`$ tends to $`\mathrm{}`$. Intuitively, in order to chose $`z_i`$ as far as possible from $`H`$, the points $`\theta _i`$ have to be chosen tranversally to $`\mathrm{\Lambda }_C`$. This transversality condition is not well defined because the limit set $`\mathrm{\Lambda }_C`$ might be highly non regular. Thus, in order to prove that such a choice is possible, we argue again by contradiction. If for any choice of a sequence $`\theta _i`$ converging to $`\theta _0`$, the distance between $`z_i`$ and $`H`$ stays bounded, then the Gromov distances $`d(\theta _i,\theta _0)`$ between $`\theta _i`$ and $`\theta _0`$ satisfy $`d(\theta _i,\mathrm{\Lambda }_C)=o\left(d(\theta _i,\theta _0)\right)`$, and therefore any tangent cone of $`\mathrm{\Lambda }_C`$ at $`\theta _0`$ would coincide with a tangent cone of $`\stackrel{~}{X}`$ at $`\theta _0`$, which is known to be topologically $`^{n1}`$. But on the other hand, the existence of a point $`x`$ such that $`\left|Jac_{n1}\left(x\right)\right|=1`$ and the fact that $`C`$ acts uniformly quasiconformally with respect to the Gromov distance on $`\stackrel{~}{X}`$ imply that the Alexandroff compactification of the above tangent cone of $`\mathrm{\Lambda }_C`$ at $`\theta _0`$ is homeomorphic to $`\mathrm{\Lambda }_C`$ which is contained in a topological sphere $`S^{n2}`$, leading to a contradiction.
From convex cocompactness of $`C`$ and the fact that the limit set of $`C`$ is a topological $`\left(n2\right)`$-dimensional sphere, there is an alternative proof of the existence of a totally geodesic $`C`$-invariant copy of the hyperbolic space $`_{}^{n1}\stackrel{~}{X}`$ which consists in observing that the topological dimension and the Hausdorff dimension of the limit set $`\mathrm{\Lambda }\left(C\right)`$ are equal to $`n2`$ and then use the following result of M. Bonk and B. Kleiner (which we quote in the riemannian manifold setting although it remains true for $`CAT\left(1\right)`$ spaces) instead of the (simpler) minimal current argument.
###### Theorem 1.7.
Let $`X`$ be a Cartan Hadamard $`n`$-dimensional manifold whose sectional curvature satisfy $`K1`$, and $`C`$ a convex cocompact discrete subgroup of isometries of $`X`$ with limit set $`\mathrm{\Lambda }_C`$. Let us assume that the topological dimension and the Hausdorff dimension (with respect to the Gromov distance on $`\stackrel{~}{X}`$) of $`\mathrm{\Lambda }_C`$ coincide and are equal to an integer $`p`$. Then, $`C`$ preseves a totally geodesic embedded copy of the real hyperbolic space $`^{p+1}`$, with $`^{p+1}=\mathrm{\Lambda }_C`$.
The authors would like to express their gratitude to Alex Lubotsky, Jean Barge, Marc Bourdon, Gilles Carron, Jean Lanne, Frédéric Paulin, Leonid Potyagailo for their interests and helpfull conversations.
## 2. Essential hypersurfaces
Let $`\mathrm{\Gamma }`$ be a discrete cocompact group of isometries of a $`n`$-dimensional Cartan-Hadamard manifold $`(\stackrel{~}{X},\stackrel{~}{g})`$ whose sectional curvature satisfies $`K_{\stackrel{~}{g}}1`$. Let us assume that the compact manifold $`X=\stackrel{~}{X}/\mathrm{\Gamma }`$ is orientable. Let us also assume that $`\mathrm{\Gamma }=A_CB`$ is an amalgamated product of its subgroups $`A`$ and $`B`$ over $`C`$.
We first reduce to the case where $`\left[\mathrm{\Gamma }:C\right]`$ is infinite.
Namely, if $`[\mathrm{\Gamma }:C]<\mathrm{}`$, then the critical exponent $`\delta \left(C\right)=\delta \left(\mathrm{\Gamma }\right)n1`$, and the equality in theorem 1.2 holds.
We then can assume that $`[\mathrm{\Gamma }:C]=\mathrm{}`$.
###### Lemma 2.1.
Let $`\mathrm{\Gamma }=A_CB`$ be as above, with $`[\mathrm{\Gamma }:C]=\mathrm{}`$. If neither $`A`$ nor $`B`$ equals $`\mathrm{\Gamma }`$, then $`H_{n1}(\stackrel{~}{X}/C,)0`$.
Proof : The Mayer-Vietoris sequence coming from the decomposition $`\mathrm{\Gamma }=A_CB`$ writes cf. , Corollary 7.7,
$$H_n(\stackrel{~}{X}/C,)H_n(\stackrel{~}{X}/A,)H_n(\stackrel{~}{X}/B,)H_n(\stackrel{~}{X}/\mathrm{\Gamma },)\mathrm{}$$
$$\mathrm{}H_{n1}(\stackrel{~}{X}/C,)\mathrm{}$$
As $`[\mathrm{\Gamma }:C]=\mathrm{}`$, $`H_n(\stackrel{~}{X}/C,)=0`$ thus, if $`H_{n1}(\stackrel{~}{X}/C,)=0`$, we deduce from the Mayer-Vietoris sequence that $`H_n(\stackrel{~}{X}/A,)H_n(\stackrel{~}{X}/B,)`$ is isomorphic to $`H_n(\stackrel{~}{X}/\mathrm{\Gamma },)`$. As $`H_{n1}(\stackrel{~}{X}/\mathrm{\Gamma },)=`$, we then deduce that either $`[\mathrm{\Gamma }:A]=\mathrm{}`$ and $`B=\mathrm{\Gamma }`$, or $`[\mathrm{\Gamma }:B]=\mathrm{}`$ and $`A=\mathrm{\Gamma }`$. $`\mathrm{}`$
In fact in the sequel of the paper we will make use of a smooth essential hypersurface $`Z`$ in $`\stackrel{~}{X}/C`$.
###### Definition 2.2.
A compact smooth orientable hypersurface $`Z`$ of an $`n`$-dimensional manifold $`Y`$ is essential in $`Y`$ if $`i_{}([Z])0`$ where $`[Z]H_{n1}(Z,)`$ denotes the fundamental class of $`Z`$ and $`i_{}:H_{n1}(Z,)H_{n1}(Y,)`$ the morphism induced by the inclusion $`i:ZY`$.
The end of this section is devoted to finding such an hypersurface $`Z`$ in $`\stackrel{~}{X}/C`$.
Let us recall a few facts about amalgamated products and their actions on trees, following . Let $`\mathrm{\Gamma }=A_CB`$ be an amalgamated products of its subgroups $`A`$ and $`B`$ over $`C`$. Then, $`\mathrm{\Gamma }`$ acts on a simplicial tree $`\stackrel{~}{T}`$ with a fundamental domain $`T\stackrel{~}{T}`$ being a segment, ie. an edge joining two vertices. Let us describe this tree $`\stackrel{~}{T}`$. There are two orbits of vertices $`\mathrm{\Gamma }v_A`$ and $`\mathrm{\Gamma }v_B`$, the stabilizer of the edge $`v_A`$ (resp. $`v_B`$) being $`A`$,( resp. $`B`$). There is one orbit of edges $`\mathrm{\Gamma }e_C`$, the stabilizer of the edge $`e_C`$ being $`C`$. The fundamental domain $`T`$ can be chosen as the edge $`e_C`$ joining the two vertices $`v_A`$ and $`v_B`$. The set of vertices adjacent to $`v_A`$, (resp. $`v_B`$), is in one to one correspondance with $`A/C`$, (resp. $`B/C`$). Note that as neither $`A`$ nor $`B`$ are equal to $`\mathrm{\Gamma }`$, then $`[A:C]1`$ and $`[B:C]1`$, therefore for an arbitrary point $`t_0`$ on the edge $`e_C`$ we see that $`\stackrel{~}{T}t_0`$ is a disjoint union of two unbounded connected components. This fact will be used later on.
Let us consider a continuous $`\mathrm{\Gamma }`$-equivariant map $`\stackrel{~}{f}:\stackrel{~}{X}T`$ where $`T`$ is the Bass-Serre tree associated to the amalgamation $`\mathrm{\Gamma }=A_CB`$. One regularizes $`\stackrel{~}{f}`$ such that it is smooth in restriction to the complementary of the inverse image of the set of vertices of $`T`$. Let $`t_0`$ a regular value of $`\stackrel{~}{f}`$ contained in that edge of $`T`$ which is fixed by the subgroup $`C`$ and define $`\stackrel{~}{Z}=\stackrel{~}{f}^1\left(t_0\right)`$. $`\stackrel{~}{Z}`$ is a smooth orientable possibly not connected hypersurface in $`\stackrel{~}{X}`$, globally $`C`$-invariant. Let us write $`Z=\stackrel{~}{Z}/C`$. We will show $`Z\stackrel{~}{X}/C`$ is compact and that one of the connected components of $`Z`$ is essential.
###### Lemma 2.3.
$`Z\stackrel{~}{X}/C`$ is compact.
Proof : Let us show that for any sequence $`z_n\stackrel{~}{Z}`$, there exists a subsequence $`z_{n_k}`$ and $`\gamma _kC`$ such that $`\gamma _kz_{n_k}`$ converges. As $`\mathrm{\Gamma }`$ is cocompact, there exists $`g_n\mathrm{\Gamma }`$ such that the set $`\left(g_nz_n\right)`$ is relatively compact. Let $`g_{n_k}z_{n_k}`$ a subsequence which converges to a point $`z\stackrel{~}{X}`$. By continuity, the sequence $`\stackrel{~}{f}\left(g_{n_k}z_{n_k}\right)`$ converges to $`\stackrel{~}{f}\left(z\right)`$, and by equivariance we get
$$g_{n_k}\stackrel{~}{f}\left(z_{n_k}\right)=g_{n_k}t_0\stackrel{~}{f}\left(z\right)$$
when $`k`$ tends to $`\mathrm{}`$. As $`\mathrm{\Gamma }`$ acts in a simplicial way on the tree $`T`$ and transitively on the set of edges, the sequence $`g_{n_k}t_0`$ is stationary, ie $`g_{n_k}t_0=t_0^{}=gt_0`$ for $`k`$ large enough. Thus $`g^1g_{n_k}=\gamma _kC`$ for $`k`$ large enough since it fixes $`t_0`$ and $`\gamma _kz_{n_k}=g^1g_{n_k}z_{n_k}`$ converges to $`g^1\left(z\right)`$ $`\mathrm{}`$
The smooth compact hypersurface $`Z`$ we constructed might be not connected. Let us write $`Z=Z_1Z_2\mathrm{}Z_k`$ where the $`Z_j`$’s are the connected components of $`Z`$. Each $`Z_j`$ is a compact smooth oriented hypersurface of $`\stackrel{~}{X}/C`$.
The aim of what follows is to prove that at least one component $`Z_i`$ of $`Z`$ is essential.
###### Lemma 2.4.
There exists $`i[1,k]`$ such that $`Z_i`$ is essential in $`\stackrel{~}{X}/C`$.
Proof : If there exists a $`Z_i`$ which doesn’t separate $`\stackrel{~}{X}/C`$ in two connected components, then $`Z_i`$ is essential in $`\stackrel{~}{X}/C`$. So we can assume that every $`Z_j`$, $`j=1,\mathrm{}k`$, does separate $`\stackrel{~}{X}/C`$ in two connected components. In that case we will show that there exists a $`Z_i`$ which separates $`\stackrel{~}{X}/C`$ in two unbounded connected components which easily implies that $`Z_i`$ is essential.
Let us denote $`U_l`$, $`l=1,2,\mathrm{},p`$, the connected components of $`\stackrel{~}{X}/C_{j=1}^kZ_j`$.
Claim : at least two components $`U_m`$, $`U_m^{}`$ are unbounded.
Assuming the claim let us finish the proof of the lemma. For each $`Z_j`$ we denote $`V_j`$, $`V_j^{}`$ the two connected components of $`\stackrel{~}{X}/CZ_j`$. Then $`U_m=W_1W_2\mathrm{}W_k`$ where for each $`j`$, $`W_j=V_j`$ or $`W_j=V_j^{}`$. In the same way, $`U_m^{}=W_1^{}W_2^{}\mathrm{}W_k^{}`$. As $`U_mU_m^{}=\mathrm{}`$, there exists $`i[1,k]`$ such that $`W_iW_i^{}=\mathrm{}`$, thus $`U_mV_i`$ and $`U_m^{}V_i^{}`$ or $`U_mV_i^{}`$ and $`U_m^{}V_i`$ so $`Z_i`$ separates $`\stackrel{~}{X}/C`$ into two unbounded components. This proves the lemma.
Let us prove the claim.
We have already noticed that $`T\left\{t_0\right\}`$ is the disjoint union of two unbounded connected components $`T_1`$ and $`T_2`$. As $`C`$ acts on $`T`$ isometrically and simplicially then $`T/C\left\{t_0\right\}=T_1/CT_2/C`$ is the disjoint union of two unbounded connected components. Let $`\overline{f}:\stackrel{~}{X}/CT/C`$ the quotient map of $`\stackrel{~}{f}`$. For each component $`U_i`$, we have
$`\overline{f}\left(U_i\right)T_1/C`$ or $`\overline{f}\left(U_i\right)T_2/C`$, thus we can conclude the claim because $`\overline{f}`$ is onto $`\mathrm{}`$
Let $`\pi :\stackrel{~}{X}\stackrel{~}{X}/C`$ be the natural projection. For any $`i=1,2,\mathrm{},k`$, let us denote $`\left\{\stackrel{~}{Z}_i^j\right\}_{jJ}`$ the set of connected components of $`\stackrel{~}{Z}_i=:\pi ^1\left(Z_i\right)`$.
For each $`i[1,k]`$, we claim that $`C`$ acts transitively on the set $`\left\{\stackrel{~}{Z}_i^j\right\}_{jJ}`$. Namely, let us consider $`\stackrel{~}{Z}_i^j`$ $`\stackrel{~}{Z}_i^j^{}`$, $`\stackrel{~}{z}\stackrel{~}{Z}_i^j`$, $`\stackrel{~}{z^{}}\stackrel{~}{Z}_i^j^{}`$, and write $`z=\pi \stackrel{~}{z}Z_i`$ and $`z^{}=\pi \stackrel{~}{z}^{}Z_i`$. Let $`\alpha `$ be a continuous path on $`Z_i`$ such that $`\alpha \left(0\right)=z`$ and $`\alpha \left(1\right)=z^{}`$, and $`\stackrel{~}{\alpha }`$ the lift of $`\alpha `$ such that $`\stackrel{~}{\alpha }\left(0\right)=\stackrel{~}{z}`$. We have $`\pi \stackrel{~}{\alpha }\left(1\right)=z^{}`$ and $`\stackrel{~}{\alpha }\left(1\right)\stackrel{~}{Z}_i^j`$ for some $`j`$, thus, there exists $`cC`$ such that $`c\left(\stackrel{~}{\alpha }\left(1\right)\right)=\stackrel{~}{z}^{}`$ and therefore $`c\stackrel{~}{Z}_i^j=\stackrel{~}{Z}_i^j^{}`$.$`\mathrm{}`$
Let us denote $`C_i^j`$ the stabilizer of $`\stackrel{~}{Z}_i^j`$, and $`Z_i^j=\stackrel{~}{Z}_i^j/C_i^j\stackrel{~}{X}/C_i^j`$. Let us write $`p:\stackrel{~}{X}/C_i^j\stackrel{~}{X}/C`$ the natural projection.
###### Lemma 2.5.
The restriction of $`p`$ to $`Z_i^j`$ is a diffeomorphism onto $`Z_i`$. In particular, $`Z_i^j`$ is compact.
proof : Let $`z`$ and $`z^{}`$ be two points in $`Z_i^j`$ such that $`p\left(z\right)=p\left(z^{}\right)`$. Let $`\stackrel{~}{z}`$ and $`\stackrel{~}{z^{}}`$ be lifts of $`z`$ and $`z^{}`$ in $`\stackrel{~}{X}`$. These two points $`\stackrel{~}{z}`$ and $`\stackrel{~}{z^{}}`$ which are in $`\stackrel{~}{Z}_i`$ actually belong to the same connected component $`\stackrel{~}{Z}_i^j`$ because for $`jj^{}`$, $`\stackrel{~}{Z}_i^j^{}/C_i^j\stackrel{~}{Z}_i^j/C_i^j=\mathrm{}`$. As $`p\left(z\right)=p\left(z^{}\right)`$, there exits $`cC`$ such that $`\stackrel{~}{z}^{}=c\stackrel{~}{z}`$, thus $`cC_i^j`$, and $`z=z^{}`$, therefore the restriction of $`p`$ to $`Z_i^j`$ is injective. The surjectivity comes from the fact that $`\pi ^1Z_i=_{jJ}\stackrel{~}{Z}_i^j`$ and $`C`$ acts transitively on the set $`\left\{\stackrel{~}{Z}_i^j\right\}_{jJ}`$. $`\mathrm{}`$
Let us consider the integer $`i[1,k]`$ as in lemma 2.4, ie. such that $`Z_i\stackrel{~}{X}/C`$ is essential, and choose $`\stackrel{~}{Z}_i^l`$ one component of $`\pi ^1\left(Z_i\right)`$.
After possibly replacing $`C_i^l`$ by an index two subgroup, we may assume that $`C_i^l`$ globally preserves each of the two connected components $`U_i^l`$ and $`V_i^l`$ of $`\stackrel{~}{X}\stackrel{~}{Z}_i^l`$.
###### Lemma 2.6.
Let $`i`$, $`l`$ and $`C_i^l`$ be chosen as above. The compact hypersurface $`Z_i^l=\stackrel{~}{Z}_i^l/C_i^l`$ is essential in $`\stackrel{~}{X}/C_i^l`$. Moreover the two connected components $`U_i^l/C_i^l`$ and $`V_i^l/C_i^l`$ of $`\stackrel{~}{X}/C_i^lZ_i^l`$ are unbounded.
Proof : Let us consider $`p:\stackrel{~}{X}/C_i^l\stackrel{~}{X}/C`$. By lemma 2.5, the restriction of $`p`$ to $`Z_i^l`$ is a diffeomorphism onto $`Z_i`$, therefore $`Z_i^l`$ is essential in $`\stackrel{~}{X}/C_i^l`$ because $`Z_i`$ is essential in $`\stackrel{~}{X}/C`$. As $`C_i^l`$ preseves $`U_i^l`$ and $`V_i^l`$, $`Z_i^l`$ separates $`\stackrel{~}{X}/C_i^l`$ into two connected components $`U_i^l/C_i^l`$ and $`V_i^l/C_i^l`$. and as $`Z_i^l`$ is essential in $`\stackrel{~}{X}/C_i^l`$, $`U_i^l/C_i^l`$ and $`V_i^l/C_i^l`$ are unbounded. $`\mathrm{}`$
In the sequel of the paper we will denote $`\stackrel{~}{Z^{}}=\stackrel{~}{Z}_i^l`$, $`C^{}=C_i^l`$ and $`Z^{}=Z_i^l=\stackrel{~}{Z}_i^l/C_i^l`$.
## 3. Isosystolic inequality
In this section we summarize facts and results due to M.Gromov, . Let $`Z`$ be a $`p`$-dimensional compact orientable manifold and $`i:ZY`$ an embedding of $`Z`$ into $`Y`$ where $`Y`$ is an aspherical space. We suppose that $`i_{}\left(\left[Z\right]\right)0`$ where $`i_{}:H_p(Z,)H_p(Y,)`$ is the morphism induced by the embedding $`ZY`$. Let us fix a riemanniann metric $`g`$ on $`Z`$. For each $`zZ`$ we consider the set $`𝒞_z`$ of those loops $`\alpha `$ at $`z`$ such that $`i\alpha `$ is homotopically non trivial in $`Y`$.
Let us define the systole of $`(Z,g,i)`$ at the point $`z`$ by
###### Definition 3.1.
$`sys_i(Z,g,z)=inf\{lengh\left(\alpha \right)`$, $`\alpha 𝒞_z\}`$
and the systole of $`(Z,g,i)`$ by
###### Definition 3.2.
$`sys_i(Z,g)=inf\{sys_i(Z,g,z)`$ $`zZ\}`$.
The following isosystolic inequality, due to M.Gromov says that the volume of any essential submanifold $`Z`$ of an aspherical space $`Y`$ relatively to any riemanniann metric on $`Z`$ is universally bounded below by it’s systole.
###### Theorem 3.3.
There exists a constant $`C_p`$ such that for each p-dimensional riemanniann manifold $`(Z,g)`$ and any embedding $`ZY`$ into an aspherical space $`Y`$ such that $`i_{}([Z])0`$ where $`i_{}:H_p(Z,)H_p(Y,)`$ is the induced morphism in homology, then $`vol_p(Z,g)C_p(sys_i(Z,g))^p`$
We will apply this volume estimates to the essential hypersurface $`i:Z\stackrel{~}{X}/C`$ that we constructed in lemma 2.6.
The following lemma is immediate.
###### Lemma 3.4.
Let $`C`$ be a discrete group acting on a simply connected manifold $`\stackrel{~}{X}`$, $`\stackrel{~}{Z}`$ a $`C`$-invariant hypersurface of $`\stackrel{~}{X}`$ and $`i:Z=\stackrel{~}{Z}/C\stackrel{~}{X}/C`$ the natural inclusion. Let $`g`$ any riemanniann metric on $`Z`$ and $`\stackrel{~}{g}`$ the lift of $`g`$ to $`\stackrel{~}{Z}`$. Then, for any $`zZ`$ we have,
$$sys_i(Z,g,z)=inf\left\{d_{\stackrel{~}{g}}(\stackrel{~}{z},\gamma \stackrel{~}{z}),\gamma C\right\}$$
where $`\stackrel{~}{z}\stackrel{~}{Z}`$ is a lift of $`zZ`$ and $`d_{\stackrel{~}{g}}`$ is the distance induced by $`\stackrel{~}{g}`$ on $`\stackrel{~}{Z}`$.
Proof : Let $`\alpha 𝒞_z`$ a loop based at $`zZ`$. As $`i\alpha `$ is an homotopically non trivial loop at $`i\left(z\right)=z`$ in $`\stackrel{~}{X}/C`$, its lift $`\stackrel{~}{i\alpha }`$ at some $`\stackrel{~}{z}\stackrel{~}{Z}`$ ends up at $`\gamma \stackrel{~}{z}`$ for some $`\gamma C`$.$`\mathrm{}`$
## 4. Volume of hypersurfaces in $`\stackrel{~}{X}/C`$
Let $`(\stackrel{~}{X},\stackrel{~}{g})`$ be a $`n`$-dimensional Cartan-Hadamard manifold whose sectional curvature satisfies $`K_{\stackrel{~}{g}}1`$ and $`C`$ a discrete group of isometries of $`(\stackrel{~}{X},\stackrel{~}{g})`$. We assume that the group $`C`$ is non elementary, namely $`C`$ fixes neither one nor two points in the geometric boundary $`\stackrel{~}{X}`$ of $`(\stackrel{~}{X},\stackrel{~}{g})`$.
The aim of this section is to construct a map $`F:\stackrel{~}{X}/C\stackrel{~}{X}/C`$ such that for any compact hypersurface $`Z`$ of $`\stackrel{~}{X}/C`$, we have
$$vol_{n1}\left(F\left(Z\right)\right)\left(\frac{\delta +1}{n1}\right)^{n1}vol_{n1}\left(Z\right)$$
where $`\delta `$ is the critical exponent of $`C`$ and $`vol_{n1}\left(Z\right)`$ stands for the $`\left(n1\right)`$-dimensional volume of the metric on $`Z`$ induced from $`g`$. For every subgroup $`C^{}C`$ and any hypersurface $`Z^{}`$ of $`\stackrel{~}{X}/C^{}`$ the lift $`F^{}:\stackrel{~}{X}/C^{}\stackrel{~}{X}/C^{}`$ of $`F`$ will also verify
$$vol_{n1}\left(F^{}\left(Z^{}\right)\right)\left(\frac{\delta +1}{n1}\right)^{n1}vol_{n1}\left(Z^{}\right).$$
In order to construct the map $`F`$ we need a few prelimiraries. We consider a finite positive Borel measure $`\mu `$ on the boundary $`\stackrel{~}{X}`$ whose support contains at least two points. Let us fix an origin $`o\stackrel{~}{X}`$ and denote $`B(x,\theta )`$ the Busemann function defined for each $`x\stackrel{~}{X}`$ and $`\theta \stackrel{~}{X}`$ by
$$B(x,\theta )=\underset{t\mathrm{}}{lim}dist(x,c\left(t\right))t$$
where $`c\left(t\right)`$ is the geodesic ray such that $`c\left(0\right)=o`$ and $`c\left(+\mathrm{}\right)=\theta `$.
Let $`𝒟_\mu :\stackrel{~}{X}`$ the function defined by
(4.1)
$$𝒟_\mu \left(y\right)=_{\stackrel{~}{X}}e^{B(y,\theta )}𝑑\mu \left(\theta \right)$$
A computation shows that
(4.2)
$$Dd𝒟_\mu \left(y\right)=_{\stackrel{~}{X}}\left(DdB(y,\theta )+DB(y,\theta )DB(y,\theta )\right)e^{B(y,\theta )}𝑑\mu \left(\theta \right).$$
When $`K_{\stackrel{~}{g}}1`$ the Rauch comparison theorem says that for every $`y\stackrel{~}{X}`$, and $`\theta \stackrel{~}{X}`$,
(4.3)
$$DdB(y,\theta )+DB(y,\theta )DB(y,\theta )\stackrel{~}{g}.$$
We then get
(4.4)
$$Dd𝒟_\mu \left(y\right)𝒟_\mu \left(y\right)\stackrel{~}{g},$$
thus $`Dd𝒟_\mu \left(y\right)`$ is positive definite and $`𝒟_\mu `$ is strictly convex.
###### Lemma 4.1.
We have $`lim_{y_k\stackrel{~}{X}}𝒟_\mu (y)=+\mathrm{}.`$
proof : Let $`y_k\stackrel{~}{X}`$ a sequence such that
(4.5)
$$\underset{k\mathrm{}}{lim}y_k=\theta _0\stackrel{~}{X}.$$
As $`supp\left(\mu \right)\left(\stackrel{~}{X}\left\{\theta _0\right\}\right)\mathrm{}`$, there exists a compact subset $`K\stackrel{~}{X}\left\{\theta _0\right\}`$ such that $`\mu \left(K\right)>0`$ thus,
(4.6)
$$_{\stackrel{~}{X}}e^{B(y_k,\theta )}𝑑\mu _Ke^{B(y_k,\theta )}𝑑\mu +\mathrm{}.$$
$`\mathrm{}`$
###### Corollary 4.2.
Let $`\mu `$ a finite borel measure on $`\stackrel{~}{X}`$ whose support contains at least two points. The function $`𝒟_\mu `$ has a unique minimum. This minimum will be denoted by $`𝒞(\mu )`$.
Let us now consider some discrete subgroup $`CIsom(\stackrel{~}{X},\stackrel{~}{g})`$. Recall that a family of Patterson measures $`\left(\mu _x\right)_{x\stackrel{~}{X}}`$ associated to $`C`$ is a set of positive finite measures $`\mu _x`$ on $`\stackrel{~}{X}`$, $`x\stackrel{~}{X}`$, such that the following holds for all $`x\stackrel{~}{X}`$, $`\gamma C`$,
(4.7)
$$\mu _{\gamma x}=\gamma _{}\mu _x$$
(4.8)
$$\mu _x=e^{\delta B(x,\theta )}\mu _o,$$
where $`o\stackrel{~}{X}`$ is a fixed origin, $`B`$ the Busemann function associated to $`o`$ and $`\delta `$ the critical exponent of $`C`$.
We assume now that $`supp\left(\mu _o\right)`$ contains at least two points and define the map $`\stackrel{~}{F}:\stackrel{~}{X}\stackrel{~}{X}`$ for $`x\stackrel{~}{X}`$ by
(4.9)
$$\stackrel{~}{F}\left(x\right)=𝒞\left(e^{B(x,\theta )}\mu _x\right).$$
Here are a few notations. For a subspace $`E`$ of $`T_x\stackrel{~}{X}`$, we will write $`Jac_E\stackrel{~}{F}\left(x\right)`$ the determinant of the matrix of the restriction of $`D\stackrel{~}{F}\left(x\right)`$ to $`E`$ with respect to orthonormal bases of $`E`$ and $`D\stackrel{~}{F}\left(x\right)E`$. For an integer $`p`$, we denote by $`Jac_p\stackrel{~}{F}\left(x\right)`$ the supremum of $`\left|Jac_E\stackrel{~}{F}\left(x\right)\right|`$ as $`E`$ runs through the set of $`p`$-dimensional subspaces of $`T_x\stackrel{~}{X}`$.
###### Lemma 4.3.
The map $`\stackrel{~}{F}`$ is smooth, homotopic to the Identity and verifies for all $`x\stackrel{~}{X}`$, $`\gamma C`$ and $`p[2,n=dim(X)]`$,
(i) $`\stackrel{~}{F}(\gamma x)=\gamma \stackrel{~}{F}(x)`$
(ii) $`|Jac_p\stackrel{~}{F}(x)|\left(\frac{(\delta +1)}{p}\right)^p`$.
Proof :
The map
$$(x,y)_{\stackrel{~}{X}}e^{B(y,\theta )B(x,\theta )}𝑑\mu _x\left(\theta \right)=_{\stackrel{~}{X}}e^{B(y,\theta )\left(\delta +1\right)B(x,\theta )}𝑑\mu _o\left(\theta \right)$$
is smooth because $`yB(y,\theta )`$ is smooth.
For all $`x`$ the map $`y_{\stackrel{~}{X}}e^{B(y,\theta )\left(\delta +1\right)B(x,\theta )}𝑑\mu _o\left(\theta \right)`$ is strictly convex by (4.4) and tends to infinity when $`y`$ tend to $`\stackrel{~}{X}`$ (cf. lemma 4.1), thus the unique minimum $`\stackrel{~}{F}\left(x\right)`$ is a smooth function. The equivariance of $`\stackrel{~}{F}`$ comes from the cocycle relation $`B(\gamma y,\gamma \theta )B(\gamma x,\gamma \theta )=B(y,\theta )B(x,\theta )`$.
For each $`x\stackrel{~}{X}`$ let $`c_x`$ be the geodesic in $`\stackrel{~}{X}`$ such that $`c_x\left(0\right)=x`$, $`c_x\left(1\right)=\stackrel{~}{F}\left(x\right)`$ and which is parametrized with constant speed. The map $`\stackrel{~}{F}_t:\stackrel{~}{X}\stackrel{~}{X}`$ defined by $`\stackrel{~}{F}_t\left(x\right)=c_x\left(t\right)`$ is a $`C`$-equivariant homotopy between $`Id_{\stackrel{~}{X}}`$ and $`\stackrel{~}{F}`$.
It remains to prove (ii).
The point $`\stackrel{~}{F}\left(x\right)`$ is characterized by
(4.10)
$$_{\stackrel{~}{X}}DB(\stackrel{~}{F}\left(x\right),\theta )e^{B(\stackrel{~}{F}\left(x\right),\theta )B(x,\theta )}𝑑\mu _x\left(\theta \right)=0.$$
In order to simplify the notations we will write $`B_{(x,\theta )}`$ instead of $`B(x,\theta )`$ and we will denote $`\nu _x`$ the measure $`e^{B(\stackrel{~}{F}\left(x\right),\theta )B(x,\theta )}\mu _x`$. We will also write $`D\stackrel{~}{F}\left(u\right)`$ instead of $`D\stackrel{~}{F}\left(x\right)\left(u\right)`$.
The differential of $`\stackrel{~}{F}`$ is characterized by the following: for $`uT_x\stackrel{~}{X}`$ and $`vT_{\stackrel{~}{F}\left(x\right)}\stackrel{~}{X}`$, one has
$$_{\stackrel{~}{X}}\left[DdB_{(\stackrel{~}{F}\left(x\right),\theta )}(D\stackrel{~}{F}\left(u\right),v)+DB_{(\stackrel{~}{F}\left(x\right),\theta )}\left(v\right)DB_{(\stackrel{~}{F}\left(x\right),\theta )}\left(D\stackrel{~}{F}\left(u\right)\right)\right]𝑑\nu _x\left(\theta \right)$$
(4.11)
$$=\left(\delta +1\right)_{\stackrel{~}{X}}DB_{(\stackrel{~}{F}\left(x\right),\theta )}\left(v\right)DB_{(x,\theta )}\left(u\right)𝑑\nu _x\left(\theta \right).$$
We define the quadratic forms $`k`$ and $`h`$ for $`vT_{\stackrel{~}{F}\left(x\right)}\stackrel{~}{X}`$ by
(4.12)
$$k(v,v)=_{\stackrel{~}{X}}\left[DdB_{(\stackrel{~}{F}\left(x\right),\theta )}(v,v)+\left(DB_{(\stackrel{~}{F}\left(x\right),\theta )}\left(v\right)\right)^2\right]𝑑\nu _x\left(\theta \right).$$
and
(4.13)
$$h(v,v)=_{\stackrel{~}{X}}DB_{(\stackrel{~}{F}\left(x\right),\theta )}\left(v\right)^2𝑑\nu _x\left(\theta \right).$$
The relation (4.11) writes, for $`uT_x\stackrel{~}{X}`$ and $`vT_{\stackrel{~}{F}\left(x\right)}\stackrel{~}{X}`$ :
(4.14)
$$k(D\stackrel{~}{F}\left(u\right),v)=\left(\delta +1\right)_{\stackrel{~}{X}}DB_{(\stackrel{~}{F}\left(x\right),\theta )}\left(v\right)DB_{(x,\theta )}\left(u\right)𝑑\nu _x\left(\theta \right).$$
We defines the quadratic form $`h^{}`$ on $`T_x\stackrel{~}{X}`$ for $`uT_x\stackrel{~}{X}`$ by
(4.15)
$$h^{}(u,u)=_{\stackrel{~}{X}}DB_{(x,\theta )}\left(u\right)^2𝑑\nu _x\left(\theta \right),$$
and one derives from (4.14)
(4.16)
$$\left|k(D\stackrel{~}{F}\left(x\right)\left(u\right),v)\right|\left(\delta +1\right)h(v,v)^{1/2}h^{}(u,u)^{1/2}.$$
One now can estimate $`Jac_p\stackrel{~}{F}\left(x\right)`$. Let $`PT_x\stackrel{~}{X}`$, $`dimP=p`$. If $`D\stackrel{~}{F}\left(P\right)`$ has dimension lower than $`p`$, then there is nothing to be proven. Let us assume that $`dimD\stackrel{~}{F}\left(P\right)=p`$. Denote by the same letters $`H^{}`$ \[resp. $`H`$ and $`K`$\] the selfadjoint operators (with respect to $`\stackrel{~}{g}`$) associated to the quadratic forms $`h^{}`$ \[resp. $`h`$, $`k`$\] restricted to $`P`$ \[resp. $`D\stackrel{~}{F}\left(P\right)`$\].
Let $`\left(v_i\right)_{i=1}^p`$ an orthonormal basis of $`D\stackrel{~}{F}\left(P\right)`$ which diagonalizes $`H`$ and $`\left(u_i\right)_{i=1}^p`$ an orthonormal basis of $`P`$ such that the matrix of $`KD\stackrel{~}{F}\left(x\right):PD\stackrel{~}{F}\left(P\right)`$ is triangular. Then,
(4.17)
$$detK.\left|Jac_P\stackrel{~}{F}\left(x\right)\right|\left(\delta +1\right)^p\left(\mathrm{\Pi }_{i=1}^ph(v_i,v_i)^{1/2}\right)\left(\mathrm{\Pi }_{i=1}^ph^{}(u_i,u_i)^{1/2}\right)$$
thus,
(4.18)
$$detK.\left|Jac_P\stackrel{~}{F}\left(x\right)\right|\left(\delta +1\right)^p\left(\frac{TraceH}{p}\right)^{p/2}\left(\frac{TraceH^{}}{p}\right)^{p/2}.$$
In these inequalities one can normalize the measures
$$\nu _x=e^{B(\stackrel{~}{F}\left(x\right),\theta )B(x,\theta )}\mu _x$$
such that their total mass equals one, which gives
(4.19)
$$traceH=\mathrm{\Sigma }_{i=1}^ph(v_i,v_i)1,$$
the last inequality coming from the fact that for all $`\theta \stackrel{~}{X}`$,
(4.20)
$$\mathrm{\Sigma }_{i=1}^pDB(\stackrel{~}{F}\left(x\right),\theta )\left(v_i\right)^2B(\stackrel{~}{F}\left(x\right),\theta )^2=1$$
and from the previous normalization.
Similarly,
(4.21)
$$traceH^{}=\mathrm{\Sigma }_{i=1}^ph^{}(u_i,u_i)1.$$
We then obtain with (4.18)
(4.22)
$$detK.\left|Jac_P\stackrel{~}{F}\left(x\right)\right|\left(\frac{\delta +1}{p}\right)^p.$$
Thanks to (4.3), we have $`detK1`$, so
(4.23)
$$\left|Jac_P\stackrel{~}{F}\left(x\right)\right|\left(\frac{\delta +1}{p}\right)^p.$$
We get (ii) by taking the supremum in $`P`$. $`\mathrm{}`$
As the map $`\stackrel{~}{F}:\stackrel{~}{X}\stackrel{~}{X}`$ is $`C`$-equivariant, then for every subgroup $`C^{}C`$, $`\stackrel{~}{F}`$ gives rise to a map $`F^{}:\stackrel{~}{X}/C^{}\stackrel{~}{X}/C^{}`$ and so does the homotopy $`\stackrel{~}{F}_t`$ between $`\stackrel{~}{F}`$ and $`Id_{\stackrel{~}{X}}`$.
###### Corollary 4.4.
The map $`F^{}:\stackrel{~}{X}/C^{}\stackrel{~}{X}/C^{}`$ is homotopic to the Identity map and verifies for all $`x\stackrel{~}{X}/C^{}`$ and $`p[2,n=dimX]`$
$$\left|JacF_p^{}\left(x\right)\right|\left(\frac{\delta +1}{p}\right)^p.$$
Let $`CIsom(\stackrel{~}{X},\stackrel{~}{g})`$ as above, ie such that the support of the Patterson- Sullivan measures contains at least two points and with critical exponent $`\delta `$. Let $`C^{}C`$ be a subgroup.
Let us consider an compact hypersurface $`Z^{}\stackrel{~}{X}/C^{}`$.
Denote $`F^k=F^{}F^{}\mathrm{}..F^{}`$ the composition of $`F^{}`$ $`k`$-times.
Let us write $`g_k=\left(F^k\right)^{}g`$, where $`g`$ is the metric on $`\stackrel{~}{X}/C^{}`$ induced by $`\stackrel{~}{g}`$. The symmetric 2-tensor $`g_k`$ may not be a riemanniann metric on $`\stackrel{~}{X}/C^{}`$ nor its resriction to $`Z^{}`$, so we have to modify it. For $`ϵ>0`$, the following symmetric 2-tensor $`g_{ϵ,k}`$ is a riemanniann metric on $`\stackrel{~}{X}/C^{}`$ and so is its restriction $`h_{ϵ,k}`$ to $`Z^{}`$.
(4.24)
$$g_{ϵ,k}=g_k+ϵ^2g.$$
###### Lemma 4.5.
Let $`h_{ϵ,k}`$ be the restriction of $`g_{ϵ,k}`$ to the hypersurface $`Z^{}`$ and $`g_Z^{}`$ the restriction of $`g`$ to $`Z^{}`$. Let $`\mathrm{\Phi }_{ϵ,k}:Z^{}`$ the density defined for all $`xZ^{}`$ by $`dv_{h_{ϵ,k}}(x)=\mathrm{\Phi }_{ϵ,k}(x)dv_{g_Z^{}}(x)`$. For any sequence $`ϵ_k`$ such that $`lim_k\mathrm{}ϵ_k=0`$, there exists a sequence $`ϵ_k^{}`$, $`lim_k\mathrm{}ϵ_k^{}=0`$, such that for all $`xZ^{}`$,
$$0<\mathrm{\Phi }_{ϵ_k^{},k}\left(x\right)\left|Jac_{n1}F^k\left(x\right)\right|+ϵ_k.$$
In particular,
$$\mathrm{\Phi }_{ϵ_k^{},k}\left(x\right)\left(\frac{\delta +1}{n1}\right)^{k\left(n1\right)}+ϵ_k$$
and
$$vol(Z^{},h_{ϵ_k^{},k})\left[\left(\frac{\delta +1}{n1}\right)^{k\left(n1\right)}+ϵ_k\right]vol(Z^{},g_Z^{}).$$
###### Corollary 4.6.
Under the above asumptions, if $`\delta <n2`$ there exists a sequence $`ϵ_k^{}`$ such that $`lim_k\mathrm{}ϵ_k^{}=0`$, and $`lim_k\mathrm{}vol(Z^{},h_{ϵ_k^{},k})=0`$.
Proof of lemma 4.5 :
Let us fix $`k`$ an integer. Let $`xZ^{}`$ and $`uT_xZ^{}`$. We have $`g_{ϵ,k}(u,u)=h_{ϵ,k}(u,u)=g(DF^k\left(x\right)\left(u\right),DF^k\left(x\right)\left(u\right))+ϵ^2g(u,u)`$ thus $`h_{ϵ,k}(u,u)=g(A_{x,ϵ}u,u)`$ where $`A_{x,ϵ}End\left(T_xZ^{}\right)`$ is the self adjoint operator $`A_x=DF^k\left(x\right)^{}DF^k\left(x\right)+ϵ^2Id`$, with $`DF^k\left(x\right)^{}`$ the adjoint of $`DF^k\left(x\right):(T_xZ^{},g\left(x\right))(DF^k\left(x\right)\left(T_xZ^{}\right),g\left(F^k\left(x\right)\right)).`$
By compactness of $`Z^{}`$ and continuity of $`A_{x,ϵ}`$, there exist $`ϵ_k^{}`$ such that
$$\mathrm{\Phi }_{ϵ_k^{},k}\left(x\right)=detA_{x,ϵ_k^{}}^{1/2}detA_{x,0}+ϵ_k,$$
thus
$$\mathrm{\Phi }_{ϵ_k^{},k}\left(x\right)\left|Jac_{n1}F^k\left(x\right)\right|+ϵ_k.$$
The lemma then follows from Corollary 4.4. $`\mathrm{}`$
## 5. Proof of Theorem 1.2
This section is devoted to the proof of the Theorem 1.2. Let $`\mathrm{\Gamma }`$ be a discrete cocompact group of isometries of a $`n`$-dimensional Cartan-Hadamard manifold $`(\stackrel{~}{X},\stackrel{~}{g})`$ whose sectional curvature satisfies $`K_{\stackrel{~}{g}}1`$. We assume that $`\mathrm{\Gamma }=A_CB`$. At the end of section 2 we constructed a subgroup $`C^{}C`$ and an orientable hypersurface $`\stackrel{~}{Z^{}}\stackrel{~}{X}`$ such that $`C^{}.\stackrel{~}{Z^{}}=\stackrel{~}{Z^{}}`$ and $`Z^{}=\stackrel{~}{Z^{}}/C^{}`$ is compact in $`\stackrel{~}{X}/C^{}`$. Moreover $`Z^{}`$ is essential in $`\stackrel{~}{X}/C^{}`$ ie $`i_{}\left(\left[Z^{}\right]\right)0`$ where $`i_{}:H_{n1}(Z^{},)H_{n1}(\stackrel{~}{X}/C^{},)`$ is the morphism induced on homology groups by the inclusion $`i:Z^{}\stackrel{~}{X}/C^{}`$ and $`\left[Z^{}\right]`$ the fundamental class of $`Z^{}`$.
5.1 Proof of the inequality
We now prove the inequality in the theorem 1.2. Let us assume that $`\delta <n2`$ and derive a contradiction. Let $`h_{ϵ_k^{},k}`$ the sequence of metric defined on $`Z^{}`$ in lemma 4.5, then by corollary 4.6 we have
(5.1)
$$\underset{k\mathrm{}}{lim}vol(Z^{},h_{ϵ_k^{},k})=0.$$
We now show that the systole of the metric $`h_{ϵ_k^{},k}`$ on $`Z^{}`$ is bounded below independently of $`k`$. Recall that the systole of $`i:Z^{}\stackrel{~}{X}/C^{}`$ at a point $`zZ^{}`$ with respect to a metric $`h_{ϵ,k}`$ can be defined by
(5.2)
$$\mathrm{sys}_\mathrm{i}(\mathrm{Z}^{},\mathrm{h}_{ϵ,\mathrm{k}},\mathrm{z})=\mathrm{inf}_{\gamma \mathrm{C}^{}}\mathrm{dist}_{(\stackrel{~}{\mathrm{Z}^{}},\stackrel{~}{\mathrm{h}}_{ϵ,\mathrm{k}})}(\stackrel{~}{\mathrm{z}},\gamma \stackrel{~}{\mathrm{z}})$$
where $`\stackrel{~}{z}`$ is any lift of $`z`$ and $`\stackrel{~}{h}_{ϵ,k}`$ the lift on $`\stackrel{~}{Z^{}}`$ of $`h_{ϵ,k}`$, (cf. lemma 3.4).
Let $`\alpha \left(t\right)`$ be a minimizing geodesic between $`\stackrel{~}{z}`$ and $`\gamma \stackrel{~}{z}`$ on $`(Z^{},\stackrel{~}{h}_{ϵ,k})`$. By definition of $`\stackrel{~}{h}_{ϵ,k}`$ we have
(5.3)
$$\mathrm{dist}_{(\stackrel{~}{\mathrm{Z}^{}},\stackrel{~}{\mathrm{h}}_{ϵ,\mathrm{k}})}(\stackrel{~}{\mathrm{z}},\gamma \stackrel{~}{\mathrm{z}})\mathrm{l}_{\stackrel{~}{\mathrm{g}}}\left(\stackrel{~}{\mathrm{F}}^\mathrm{k}\alpha \right)$$
where $`l_{\stackrel{~}{g}}`$ stands for the lengh with respect to $`\stackrel{~}{g}`$ on $`\stackrel{~}{X}`$.
We get
(5.4)
$$\mathrm{dist}_{(\stackrel{~}{\mathrm{Z}^{}},\stackrel{~}{\mathrm{h}}_{ϵ,\mathrm{k}})}(\stackrel{~}{\mathrm{z}},\gamma \stackrel{~}{\mathrm{z}})\mathrm{dist}_{(\stackrel{~}{\mathrm{X}},\stackrel{~}{\mathrm{g}})}(\stackrel{~}{\mathrm{F}}^\mathrm{k}\left(\stackrel{~}{\mathrm{z}}\right),\gamma \stackrel{~}{\mathrm{F}}^\mathrm{k}\left(\stackrel{~}{\mathrm{z}}\right))\rho $$
where $`\rho `$ is the injectivity radius of $`\stackrel{~}{X}/C^{}`$.
We then have
(5.5)
$$\mathrm{sys}_\mathrm{i}(\mathrm{Z}^{},\mathrm{h}_{ϵ,\mathrm{k}})\rho ,$$
and thanks to the Theorem 3.3 we obtain
(5.6)
$$\mathrm{vol}(\mathrm{Z}^{},\mathrm{h}_{ϵ_\mathrm{k}^{},\mathrm{k}})\mathrm{C}_\mathrm{n}\rho ^{\mathrm{n}1}$$
which contradicts (5.1).
$`\mathrm{}`$
5.2 Proof of the equality case
There will be several steps.
Step 1: The limit set of $`C`$ is contained in a topological equator.
Step 2: The weak tangent to $`\stackrel{~}{X}`$ and $`\mathrm{\Lambda }_C^{}`$.
Step 3: The limit set $`\mathrm{\Lambda }_C^{}`$ of $`C^{}`$ and the limit set $`\mathrm{\Lambda }_C`$ of $`C`$ are equal to a topological equator.
Step 4: $`C^{}`$ and $`C`$ are convex cocompact.
Step 5: $`C`$ preserves a copy of the real hyperbolic space $`_{}^{n1}`$ totally geodesically embedded in $`\stackrel{~}{X}`$.
Step 6: Conclusion
Step 1: The limit set $`\mathrm{\Lambda }_C`$ of $`C`$ is contained in a topological equator.
Let $`x\stackrel{~}{X}`$ and $`ET_x\stackrel{~}{X}`$ a codimension one subspace. For each $`uT_x\stackrel{~}{X}`$, $`\stackrel{~}{g}(u,u)=1`$, one considers the geodesic $`c_u`$ defined by $`c_u\left(0\right)=x`$ and $`\dot{c_u}\left(0\right)=u`$. We define the equator $`E\left(\mathrm{}\right)`$ associated to $`E`$ as the subset of $`\stackrel{~}{X}`$
(5.7)
$$E\left(\mathrm{}\right)=\left\{c_u\left(+\mathrm{}\right)/uE\right\}$$
Our goal is to prove the existence of a point $`x\stackrel{~}{X}`$ and an hyperplane $`ET_x\stackrel{~}{X}`$ such that the limit set $`\mathrm{\Lambda }_C`$ satisfies $`\mathrm{\Lambda }_CE\left(\mathrm{}\right)`$.
Recall that $`C^{}C`$ globally preserves an hypersurface $`\stackrel{~}{Z^{}}`$ such that $`\stackrel{~}{Z^{}}/C^{}\stackrel{~}{X}/C^{}`$ is compact and essential.
Let us also recall that we have constructed a $`C`$-equivariant map $`\stackrel{~}{F}:\stackrel{~}{X}\stackrel{~}{X}`$ such that, for all $`x\stackrel{~}{X}`$,
(5.8)
$$\left|Jac_{n1}\stackrel{~}{F}\left(x\right)\right|\left(\frac{\delta +1}{n1}\right)^{n1}$$
where the critical exponent $`\delta `$ of $`C`$ satisfies $`\delta =n2`$, thus
(5.9)
$$\left|Jac_{n1}\stackrel{~}{F}\left(x\right)\right|1.$$
The step 1 follows from the two following Propositions.
###### Proposition 5.1.
Let $`x\stackrel{~}{X}`$ such that $`|Jac_{n1}\stackrel{~}{F}(x)|=1`$. Then there exists $`ET_x\stackrel{~}{X}`$ such that the limit set $`\mathrm{\Lambda }_C`$ satisfies $`\mathrm{\Lambda }_CE(\mathrm{})`$. Moreover, $`\stackrel{~}{F}(x)=x`$ and $`D\stackrel{~}{F}(x)`$ is the orthogonal projector onto $`E`$.
###### Proposition 5.2.
There exists $`x\stackrel{~}{X}`$ such that $`|Jac_{n1}\stackrel{~}{F}(x)|=1`$.
Proof of Proposition 5.1
Let $`x\stackrel{~}{X}`$ such that $`\left|Jac_{n1}\stackrel{~}{F}\left(x\right)\right|=1`$. By definition we have a subspace $`ET_x\stackrel{~}{X}`$ such that $`\left|Jac_E\stackrel{~}{F}\left(x\right)\right|=1`$. By (4.18) and $`\mathrm{detK}1`$ we have,
$$\left|Jac_E\stackrel{~}{F}\left(x\right)\right|\left(n1\right)^{n1}\left(\frac{traceH}{n1}\right)^{\frac{n1}{2}}\left(\frac{traceH^{}}{n1}\right)^{\frac{n1}{2}}$$
(5.10)
$$\left(n1\right)^{n1}\left(\frac{1}{n1}\right)^{n1}.$$
In particular as $`\left|Jac_E\stackrel{~}{F}\left(x\right)\right|=1`$, we have equality in the inequalities (5.10), thus, $`traceH=trace\left(h\right)=1`$, and
(5.11)
$$H=\frac{1}{n1}Id_{D\stackrel{~}{F}\left(x\right)\left(E\right)}.$$
Let us recall that the quadratic form $`h`$ is defined by
$$h(v,v)=_{\stackrel{~}{X}}DB(\stackrel{~}{F}\left(x\right),\theta )\left(v\right)^2e^{B(\stackrel{~}{F}\left(x\right),\theta )B(x,\theta )}𝑑\mu _x\left(\theta \right)$$
where $`\mu _x`$ is the Patterson-Sullivan measure of $`C`$ normalized by
(5.12)
$$_{\stackrel{~}{X}}e^{B(\stackrel{~}{F}\left(x\right),\theta )B(x,\theta )}𝑑\mu _x\left(\theta \right)=1.$$
We then have
$$1=trace\left(h\right)=traceH=\mathrm{\Sigma }_{i=1}^{n1}h(v_i,v_i)=$$
$$=_{\stackrel{~}{X}}\mathrm{\Sigma }_{i=1}^{n1}DB(\stackrel{~}{F}\left(x\right),\theta )\left(v_i\right)^2e^{B(\stackrel{~}{F}\left(x\right),\theta )B(x,\theta )}𝑑\mu _x\left(\theta \right)$$
$$_{\stackrel{~}{X}}DB(\stackrel{~}{F}\left(x\right),\theta )^2e^{B(\stackrel{~}{F}\left(x\right),\theta )B(x,\theta )}𝑑\mu _x\left(\theta \right)1,$$
because $`\mathrm{\Sigma }_{i=1}^{n1}DB(\stackrel{~}{F}\left(x\right),\theta )\left(v_i\right)^2DB(\stackrel{~}{F}\left(x\right),\theta )^2=1`$
for all $`\theta \stackrel{~}{X}`$.
Therefore for $`\mu _x`$-almost all $`\theta supp\left(e^{B(\stackrel{~}{F}\left(x\right),\theta )B(x,\theta )}d\mu _x\left(\theta \right)\right)=supp\left(\mu _x\right)`$, we have
(5.13)
$$\mathrm{\Sigma }_{i=1}^{n1}DB(\stackrel{~}{F}\left(x\right),\theta )\left(v_i\right)^2=DB(\stackrel{~}{F}\left(x\right),\theta )^2=1.$$
In (5.12), $`\mathrm{\Sigma }_{i=1}^{n1}DB(\stackrel{~}{F}\left(x\right),\theta )\left(v_i\right)^2`$ represents the square of the norm of the projection of $`B(\stackrel{~}{F}\left(x\right),\theta )`$ on $`E`$.
By continuity of $`B(x,\theta )`$ in $`\theta `$ one then gets $`\mathrm{\Lambda }_C=supp\left(\mu _x\right)E\left(\mathrm{}\right)`$.
Let us now prove that $`\stackrel{~}{F}\left(x\right)=x`$. When $`Jac\stackrel{~}{F}_E\left(x\right)=1`$, we have equality in the Cauchy-Schwarz inequality (4.16), therefore for each $`i=1,\mathrm{}n1`$ and $`\theta \mathrm{\Lambda }_C`$ we get $`DB(\stackrel{~}{F}\left(x\right),\theta )\left(v_i\right)=DB(x,\theta )\left(u_i\right)`$. Therefore we deduces from (5.13) that $`B(x,\theta )=\mathrm{\Sigma }_{i=1}^{i=n1}DB(x,\theta )\left(u_i\right)u_i`$, which imply with (4.10) that
(5.14)
$$_{\stackrel{~}{X}}DB(x,\theta )e^{B(\stackrel{~}{F}\left(x\right),\theta )B(x,\theta )}𝑑\mu _x\left(\theta \right)=0.$$
On the other hand, as $`_{\stackrel{~}{X}}DB(\stackrel{~}{F}\left(x\right),\theta )e^{B(\stackrel{~}{F}\left(x\right),\theta )B(x,\theta )}𝑑\mu _x\left(\theta \right)=0`$ and $`H=\frac{1}{n1}Id_{D\stackrel{~}{F}\left(x\right)\left(E\right)}`$, the support of $`\mu _x`$ cannot be just a pair of points, therefore the barycenter of the measure $`e^{B(\stackrel{~}{F}\left(x\right),\theta )B(x,\theta )}\mu _x`$ defined in is well defined and characterized as the point $`z\stackrel{~}{X}`$ such that
$$_{\stackrel{~}{X}}DB(z,\theta )e^{B(\stackrel{~}{F}\left(x\right),\theta )B(x,\theta )}𝑑\mu _x\left(\theta \right)=0,$$
thus (5.14) and (4.10) imply $`x=\stackrel{~}{F}\left(x\right)`$. $`\mathrm{}`$
Proof of Proposition 5.2 :
If we knew that there exists a minimizing hypersurface $`Z_0`$ in the homology class of $`Z`$, then every points $`xZ_0`$ would verify $`\left|Jac_{n1}\stackrel{~}{F}\left(x\right)\right|=1`$. We unfortunately don’t know if there exists such a minimizing hypersurface nor a minimizing current in the homology class of $`Z`$. Instead we will consider an $`L^2\left(\stackrel{~}{X}/C^{}\right)`$ harmonic $`\left(n1\right)`$-form dual to the homology class of $`Z`$.
We need the following lemmas in order to prove the existence of such a dual form.
Let $`\lambda _1\left(\stackrel{~}{X}/C^{}\right)`$ be the bottom of the spectrum of the Laplacian on $`(\stackrel{~}{X},\stackrel{~}{g})`$, ie.
(5.15)
$$\lambda _1\left(\stackrel{~}{X}/C^{}\right)=inf_{uC_0^{\mathrm{}}\left(\stackrel{~}{X}/C^{}\right)}\left\{\frac{_{\stackrel{~}{X}/C^{}}\left|du\right|^2}{_{\stackrel{~}{X}/C^{}}u^2}\right\}.$$
###### Lemma 5.3.
Let $`CIsom(\stackrel{~}{X},\stackrel{~}{g})`$ a discrete group of isometries with critical exponent $`\delta =n2`$ where $`(\stackrel{~}{X},\stackrel{~}{g})`$ is an n-dimensional Cartan-Hadamard manifold of sectional curvature $`K_{\stackrel{~}{g}}1`$. Then for any subgroup $`C^{}C`$ we have $`\lambda _1(\stackrel{~}{X}/C^{})n2`$.
Proof :
Thanks to a theorem of Barta, cf. Theorem 2.1, the lemma boils down to finding a positive function $`c:\stackrel{~}{X}/C^{}_+`$ such that $`\mathrm{\Delta }c\left(x\right)\left(n2\right)c\left(x\right)`$. Here, the laplacian $`\mathrm{\Delta }`$ is the positive operator ie. $`\mathrm{\Delta }c=traceDdc`$. We consider the smooth function $`\stackrel{~}{c}:\stackrel{~}{X}_+`$ defined by $`\stackrel{~}{c}\left(x\right)=\mu _x\left(\stackrel{~}{X}\right)`$ where $`\left\{\mu _x\right\}_{x\stackrel{~}{X}}`$ is a family of Patterson-Sullivan measure of $`C`$. The function $`\stackrel{~}{c}`$ is $`C`$-equivariant therefore it defines a map $`c:\stackrel{~}{X}/C^{}_+`$ for any subgroup $`C^{}C`$. Let us show
(5.16)
$$\mathrm{\Delta }c\left(x\right)\delta \left(n1\delta \right)c\left(x\right)=n2.$$
We have
$$\stackrel{~}{c}\left(x\right)=_{\stackrel{~}{X}}e^{\delta B(x,\theta )}𝑑\mu _o\left(\theta \right)$$
therefore
$$\mathrm{\Delta }\stackrel{~}{c}\left(x\right)=_{\stackrel{~}{X}}\left[\delta \mathrm{\Delta }B(x,\theta )\delta ^2\right]𝑑\mu _x\left(\theta \right).$$
The sectional curvature $`K_{\stackrel{~}{g}}`$ of $`(\stackrel{~}{X},\stackrel{~}{g})`$ satisfies $`K_{\stackrel{~}{g}}1`$ we thus have $`\mathrm{\Delta }B(x,\theta )n1`$ and as $`\delta =n2`$ we get
$$\mathrm{\Delta }\stackrel{~}{c}\left(x\right)\left[\delta \left(n1\right)\delta ^2\right]\stackrel{~}{c}\left(x\right)=\left(n2\right)\stackrel{~}{c}\left(x\right).$$
$`\mathrm{}`$
The following lemma is due to G.Carron and E.Pedon, . For a complete riemannian manifold $`Y`$, we denote $`H_c^1(Y,)`$ the first cohomology group generated by diferential forms with compact support.
###### Lemma 5.4 (, Lemme 5.1).
Let $`Y`$ be a complete riemannian manifold all ends of whose having infinite volume and such that $`\lambda _1(Y)>0`$, then the natural morphism
$$H_c^1(Y,)H_{L^2}^1(Y,)$$
is injective. In particular any $`\alpha H_c^1(Y,)`$ admits a representative $`\overline{\alpha }`$ which is in $`L^2(Y,)`$.
###### Corollary 5.5.
Let $`C^{}`$ be as above and assume that there exists a compact essential hypersurface $`Z^{}\stackrel{~}{X}/C^{}`$. Then there exists an harmonic $`n1`$-form $`\omega `$ in $`L^2(\stackrel{~}{X}/C^{})`$ such that $`_Z^{}\omega 0`$.
Proof :
Let $`\alpha H_c^1(\stackrel{~}{X}/C^{},)`$ a Poincaré dual of $`\left[Z^{}\right]H_{n1}(\stackrel{~}{X}/C^{},)`$. By definition of $`\alpha `$, for any $`\beta H^{n1}(\stackrel{~}{X}/C^{},)`$, one has
(5.17)
$$_Z^{}\beta =_{\stackrel{~}{X}/C^{}}\beta \alpha ,$$
( p.51, note that $`\stackrel{~}{X}/C^{}`$ has a ” finite good cover”).
After Lemma 5.4 , $`\alpha `$ admits a non trivial harmonic representative $`\overline{\alpha }`$ in $`L^2\left(\stackrel{~}{X}/C^{}\right)`$. (In order to apply the Lemma 5.4, one has to check that all ends of $`\stackrel{~}{X}/C^{}`$ have infinite volume, ie for a compact $`K\stackrel{~}{X}/C^{}`$ each unbounded connected component of $`\stackrel{~}{X}/C^{}K`$ has infinite volume: this comes from the fact that the injectivity radius of $`\stackrel{~}{X}/C^{}`$ is bounded below by the injectivity radius of $`X=\stackrel{~}{X}/\mathrm{\Gamma }`$ and the sectional curvature bounded above by $`1`$.) The $`\left(n1\right)`$-harmonic form $`\omega =\overline{\alpha }`$, where $``$ is the Hodge operator, is in $`L^2\left(\stackrel{~}{X}/C^{}\right)`$ and verifies after (5.17)
(5.18)
$$_Z^{}\omega =_{\stackrel{~}{X}/C^{}}\omega \overline{\alpha }=_{\stackrel{~}{X}/C^{}}\omega \omega =\left|\right|\omega ||_{L^2\left(\stackrel{~}{X}/C^{}\right)}^20.$$
$`\mathrm{}`$
We can now prove the proposition 5.2. Let us briefly describe the idea. We consider the iterates $`F^k`$ of $`F^{}:\stackrel{~}{X}/C^{}\stackrel{~}{X}/C^{}`$. As $`F^{}`$ is homotopic to the identity map, $`F^k\left(Z^{}\right)`$ is homologous to $`Z^{}`$ and if $`\omega `$ is the harmonic form of the corollary 5.5 we have
(5.19)
$$_Z^{}\left(F^k\right)^{}\left(\omega \right)=_Z^{}\omega =a0.$$
We don’t know if $`F^k\left(Z^{}\right)`$ converges or stays in a compact subset of $`\stackrel{~}{X}/C^{}`$ but we will show that $`F^k\left(Z^{}\right)`$ cannot entirely diverge in $`\stackrel{~}{X}/C^{}`$ and that there exists a $`z^{}Z^{}`$ such that $`F^k\left(z^{}\right)`$ subconverges to a point $`x\stackrel{~}{X}/C^{}`$ with $`\left|Jac_{n1}F^{}\left(x\right)\right|=1`$.
¿From (5.19) one gets
(5.20)
$$0<\left|a\right|=\left|_Z^{}\left(F^k\right)^{}\left(\omega \right)\right|_Z^{}\left|Jac_{TZ^{}}F^k\left(z\right)\right|.\left|\right|\omega \left(F^k\left(z\right)\right||dz$$
where
$$\left|Jac_{TZ^{}}F^k\left(z\right)\right|=DF^k\left(z\right)\left(u_1\right)DF^k\left(z\right)\left(u_2\right)\mathrm{}DF^k\left(z\right)\left(u_{n1}\right)$$
and $`(u_1,..,u_{n1})`$ is an orthonormal basis of $`T_z\left(Z^{}\right)`$.
Let us define $`=\{zZ^{},|JacF_{TZ^{}}^k\left(z\right)|`$ does not converge to 0 $`\}`$.
For $`zZ^{}`$ we define the sequence $`z_k`$ by $`z_0=z`$ and $`z_k=F^{}\left(z_{k1}\right)=F^k\left(z\right)\stackrel{~}{X}/C^{}`$.
###### Lemma 5.6.
There exists $`z`$ and a subsequence $`z_{k_j}`$ such that $`z_{k_j}`$ converges to a point $`x\stackrel{~}{X}/C^{}`$ with $`|Jac_{n1}(x)|=1.`$
Proof :
We first remark that $`lim_x\mathrm{}\omega \left(x\right)=0`$. This follows the following facts: $`\omega `$ is harmonic, $`\omega L^2\left(\stackrel{~}{X}/C^{}\right)`$ and the injectivity radius of $`\stackrel{~}{X}/C^{}`$ is bounded below by a positive constant.
Let us assume that for all $`z`$ the sequence $`z_k`$ diverges in $`\stackrel{~}{X}/C^{}`$. Then we have, for all $`z`$,
(5.21)
$$\omega \left(z_k\right)=\omega \left(F^k\left(z\right)\right)0$$
whenever $`k`$ tends to $`\mathrm{}`$ because of the previous remark.
On the other hand, as $`\omega \left(F^k\left(z\right)\right)C`$ and $`\left|Jac_{n1}F^k\right|1`$, it follows from (5.20)
$$\underset{k\mathrm{}}{lim}\left|_Z^{}\left(F^k\right)^{}\left(\omega \right)\right|$$
$$\underset{k\mathrm{}}{lim}\left[_{}\right|\left|\omega \left(F^k\left(z\right)\right)\right||dz+C_Z^{}|Jac_Z^{}F^k\left(z\right)|dz=0$$
which contradicts our assumption.
Thus there exists a point $`zZ^{}`$ such that
(5.22)
$$\left|Jac_Z^{}F^k\left(z\right)\right|\alpha 0$$
and such that there exists a subsequence $`z_{k_j}=F^{k_j}\left(z\right)`$ with
(5.23)
$$lim_j\mathrm{}z_{k_j}=x\stackrel{~}{X}/C^{}.$$
The property (5.22) comes from the fact that the sequence $`\left|Jac_Z^{}F^k\left(z\right)\right|`$ doesn’t tend to zero and is decreasing (because $`\left|Jac_{n1}F^{}\right|1`$).
Let us define
$$E_0=T_zZ^{},E_1=DF^{}\left(z\right)\left(E_0\right)$$
and
$$E_k=DF^{}\left(z_{k1}\right)\left(E_{k1}\right)T_{z_k}\left(\stackrel{~}{X}/C^{}\right).$$
As $`z_{k_j}x`$ we can assume, after extracting again a subsequence, that $`E_{k_j}ET_x\left(\stackrel{~}{X}/C^{}\right)`$. On the other hand we also have
(5.24)
$$\left|Jac_Z^{}F^k\left(z\right)\right|=\left|Jac_{E_{k1}}F^{}\left(z_{k1}\right)\right|\left|Jac_{E_{k2}}F^{}\left(z_{k2}\right)\right|\mathrm{}\left|Jac_{E_0}F^{}\left(z\right)\right|$$
We know that $`\left|Jac_{E_k}F^{}\left(z_k\right)\right|=1ϵ_k`$ where $`0ϵ_k<1`$. As $`z`$, we have
$$\underset{k\mathrm{}}{lim}\pi _{j=1}^k\left(1ϵ_j\right)=\alpha >0$$
therefore $`lim_k\mathrm{}ϵ_k=0`$ and by continuity we have $`\left|Jac_EF^{}\left(x\right)\right|=1.`$ $`\mathrm{}`$
Now we can finish the proof of the step 1. Let consider a lift of $`x`$ in $`\stackrel{~}{X}`$ and $`E`$ in $`T\stackrel{~}{X}`$ that we again call $`x`$ and $`E`$. Then we have $`\left|Jac_E\stackrel{~}{F^{}}\left(x\right)\right|=1`$. $`\mathrm{}`$
Let us remark that corollary 4.4 and (5.19) give another proof of the inequality $`\delta n2`$, which does not use the isosystolic inequality, ie. Theorem 3.3.
Step 2 : The weak tangent of $`\stackrel{~}{X}`$ and $`\mathrm{\Lambda }_C`$
We first recall the definition of the Gromov-distance on $`\stackrel{~}{X}`$.
For two arbitrary points $`\theta `$ and $`\theta ^{}`$ in $`\stackrel{~}{X}`$ let us define
(5.25)
$$l(\theta ,\theta ^{})=inf\left\{t>0/dist(\alpha _\theta \left(t\right),\alpha _\theta ^{}\left(t\right))=1\right\}$$
and
(5.26)
$$d(\theta ,\theta ^{})=e^{l(\theta ,\theta ^{})}$$
then $`d`$ is a distance on $`\stackrel{~}{X}`$.
We now recall a few definitions following . A complete metric space $`(S,\overline{d})`$ is a weak tangent of a metric space $`(Z,d)`$ if there exist a point $`0S`$, a sequence of points $`z_kZ`$ and a sequence of positive real numbers $`\lambda _k\mathrm{}`$ such that the sequence of pointed metric spaces $`(Z,\lambda _kd,z_k)`$ converges in the pointed Gromov-Hausdorff topology to $`(S,\overline{d},0)`$ where $`(Z,\lambda _kd)`$ stands for the set $`Z`$ endowed with the rescaled metric $`\lambda _kd`$.
Let us recall that the sequence of metric spaces $`(Z_k,d_k,z_k)`$ converges to $`(S,\overline{d},0)`$ in the pointed Gromov-Hausdorff topology if the following conditions hold, (cf. (B-B-I), definition 8.1.1).
###### Definition 5.7.
We say that the sequence of metric spaces $`(Z_k,d_k,z_k)`$ converges to $`(S,\overline{d},o)`$ if for any $`R>0`$, $`ϵ>0`$ there exists $`k_0`$ such that for any $`kk_0`$ there exists a (non necessary continuous) map $`f:B(z_k,R)S`$ such that
(i) $`f\left(z_k\right)=0`$,
for any two points $`x`$ and $`y`$ in $`B(z_k,R)`$,
(ii) $`\left|\overline{d}(f\left(x\right),f\left(y\right))d_k(x,y)\right|ϵ`$,
and
(iii) the $`ϵ`$-neighborhood of the set $`f(B(z_k,R)`$ contains $`B(0,Rϵ)`$.
In the previous definition, $`B(z_k,R)`$ stands for the ball of radius $`R`$ centered at the point $`z_k`$ in $`(Z_k,d_k)`$.
For a metric space $`(Z,d)`$ we will denote $`WT(Z,d)`$ the set of weak tangents of $`(Z,d)`$.
Let $`\mathrm{\Gamma }`$ a cocompact group of isometries of $`(\stackrel{~}{X},\stackrel{~}{g})`$ a $`n`$-dimensional Cartan Hadamard manifold of sectional curvature $`K_{\stackrel{~}{g}}1`$ and $`C^{}`$ a subgroup of $`\mathrm{\Gamma }`$. The limit set $`\mathrm{\Lambda }_\mathrm{\Gamma }`$ of $`\mathrm{\Gamma }`$ is the full boundary $`\stackrel{~}{X}`$ namely a topological $`\left(n1\right)`$dimensional sphere $`S^{n1}`$. We endow $`\stackrel{~}{X}`$ with the Gromov distance $`d`$ defined in (5.26). In Lemma 5.2, M.Bonk and B.Kleiner show among other properties the following
###### Lemma 5.8.
For any weak tangent space $`(S,\overline{d})`$ in $`WT((\stackrel{~}{X},d))`$, $`S`$ is homeomorphic to $`\stackrel{~}{X}`$ less a point, thus to $`^{n1}`$.
In fact the crucial asumption in the above lemma, coming from the cocompacness of $`\mathrm{\Gamma }`$, is the property that any triple of points in $`\stackrel{~}{X}`$ can be uniformly separated by an element of $`\mathrm{\Gamma }`$, ie. there is $`\delta >0`$ such that for any three points $`\theta _1,\theta _2,\theta _3\stackrel{~}{X}`$ there exists a $`\gamma \mathrm{\Gamma }`$ such that $`d(\gamma \theta _i,\gamma \theta _j)\delta `$ for all $`1ij3`$. Following the argument of M.Bonk and B.Kleiner one can show that if $`C^{}`$ is a subgroup of $`\mathrm{\Gamma }`$ such that one weak tangent $`(S,\overline{d})`$ of $`(\mathrm{\Lambda }_C^{},d)`$ is in $`WT(\stackrel{~}{X},d)`$ and enough triples of points of $`\mathrm{\Lambda }_C^{}`$ can be uniformly separated by elements of $`C^{}`$, then $`S`$ is homeomorphic to $`\mathrm{\Lambda }_C^{}`$ less a point. In particular, $`\mathrm{\Lambda }_C^{}`$ is homeomorphic to $`\stackrel{~}{X}`$.
###### Lemma 5.9.
Let $`\stackrel{~}{X}`$ be a closed $`C^{}`$-invariant set and $`\theta _0`$. We assume that there exist a sequence of positive real numbers $`\lambda _k\mathrm{}`$ such that the sequence of pointed metric spaces $`(,\lambda _kd,\theta _0)`$ converges in the pointed Gromov-Hausdorff topology to $`(S,\overline{d},0)`$ where $`(S,\overline{d},0)`$ is a weak tangent of $`(\stackrel{~}{X},d)`$. We also assume that there exist positive constants $`C`$ and $`\delta `$, a sequence of points $`\theta _0^k=\theta _0,\theta _1^k,\theta _2^k`$ and a sequence of elements $`\gamma _kC^{}`$ such that $`C^1\lambda _k^1d(\theta _i^k,\theta _j^k)C\lambda _k^1`$ and $`d(\gamma _k\theta _i^k,\gamma _k\theta _j^k)\delta `$ for all $`0ij2`$. Then, $`S`$ is homeomorphic to $``$ less a point. In particular $``$ is homeomorphic to $`\stackrel{~}{X}`$.
The proof of this lemma is postponed in the Appendix.
Step 3 : The limit set $`\mathrm{\Lambda }_C^{}`$ of $`C^{}`$ and the limit set $`\mathrm{\Lambda }_C`$ of $`C`$ are equal to a topological equator.
We have shown in step 1 that $`\mathrm{\Lambda }_C`$ is a subset of some topological equator $`E\left(\mathrm{}\right)`$.
Let $`o\stackrel{~}{X}`$ and $`ET_o\stackrel{~}{X}`$ be such that $`\left|Jac_E\stackrel{~}{F}\left(o\right)\right|=1`$ and $`E\left(\mathrm{}\right)`$ is the equator associated to $`E`$.
Recall that there exists a subgroup $`C^{}`$ of $`C`$ which globally preserves an hypersurface $`\stackrel{~}{Z^{}}\stackrel{~}{X}`$ and that $`\stackrel{~}{Z^{}}/C^{}\stackrel{~}{X}/C^{}`$ is compact. Furthermore $`\stackrel{~}{Z^{}}`$ separates $`\stackrel{~}{X}`$ into two connected components $`\stackrel{~}{U}`$ et $`\stackrel{~}{V}`$. We can assume that $`\stackrel{~}{U}`$ and $`\stackrel{~}{V}`$ are globally invariant by $`C^{}`$ after having replaced $`C^{}`$ by an index $`2`$ subgroup.
The limit set $`\mathrm{\Lambda }_C^{}`$ of $`C^{}`$ is contained in $`\mathrm{\Lambda }_C`$, therefore $`\mathrm{\Lambda }_C^{}E\left(\mathrm{}\right)`$.
We will show that $`\mathrm{\Lambda }_C^{}=E\left(\mathrm{}\right)`$.
For any subset $`W\stackrel{~}{X}`$ we define the boundary at infinity $`W`$ of $`W`$ by
(5.27)
$$W=Cl\left(W\right)\stackrel{~}{X}$$
where $`Cl\left(W\right)`$ stands for the closure of $`W`$ in $`\stackrel{~}{X}\stackrel{~}{X}`$.
As $`\stackrel{~}{Z^{}}/C^{}`$ is compact, $`\stackrel{~}{Z^{}}`$ is at bounded distance of the orbit $`C^{}z`$ of some point $`z`$ in $`\stackrel{~}{Z^{}}`$, thus
(5.28)
$$Cl\left(\stackrel{~}{Z^{}}\right)\stackrel{~}{X}=\mathrm{\Lambda }_C^{}$$
By definition we have $`\mathrm{\Lambda }_C^{}\stackrel{~}{U}`$ and $`\mathrm{\Lambda }_C^{}\stackrel{~}{V}`$.
###### Lemma 5.10.
Let us assume that $`\mathrm{\Lambda }_C^{}E(\mathrm{})`$, then either $`\stackrel{~}{U}=\mathrm{\Lambda }_C^{}`$ or $`\stackrel{~}{V}=\mathrm{\Lambda }_C^{}`$.
Proof: We know that $`\mathrm{\Lambda }_C^{}\stackrel{~}{U}`$ and $`\mathrm{\Lambda }_C^{}\stackrel{~}{V}`$.
Let us assume that the conclusion of the lemma is not true, so there is $`\zeta \stackrel{~}{U}\mathrm{\Lambda }_C^{}`$ and $`\theta \stackrel{~}{V}\mathrm{\Lambda }_C^{}`$.
As $`\mathrm{\Lambda }_C^{}E\left(\mathrm{}\right)`$, $`\mathrm{\Lambda }_C^{}E\left(\mathrm{}\right)`$ and $`\stackrel{~}{X}`$ is a sphere, any two points of $`\stackrel{~}{X}\mathrm{\Lambda }_C^{}`$ can be joined by a continuous path contained in $`\stackrel{~}{X}\mathrm{\Lambda }_C^{}`$ and so does $`\zeta `$ and $`\theta `$, joined by such a path $`\alpha `$.
The set $`Z^{}\mathrm{\Lambda }_C^{}`$ is a closed subset of $`\stackrel{~}{X}\stackrel{~}{X}`$ thus there is an open connected neighborhood $`W`$ of $`\alpha `$ in $`\stackrel{~}{X}\stackrel{~}{X}`$ contained in the complementary of $`Z^{}\mathrm{\Lambda }_C^{}`$.
As $`\zeta `$ and $`\theta `$ can be approximated by points in $`\stackrel{~}{U}`$ and $`\stackrel{~}{V}`$ respectively there exist points $`x\stackrel{~}{U}W`$ and $`y\stackrel{~}{V}W`$ that can be joined by a continuous path by connectedness of $`W`$, which leads to a contradiction. $`\mathrm{}`$
###### Remark 5.11.
In fact, we are going to show that under the assumption $`\delta (C^{})=n2`$, it is impossible to have $`\stackrel{~}{U}=\mathrm{\Lambda }_C^{}`$ or $`\stackrel{~}{V}=\mathrm{\Lambda }_C^{}`$.
For any $`x\stackrel{~}{X}`$ and $`\theta \stackrel{~}{X}`$ let us denote $`HB(x,\theta )`$ the open horoball centered at $`\theta `$ and passing through $`x`$.
###### Lemma 5.12.
Let us assume that $`\stackrel{~}{U}=\mathrm{\Lambda }_C^{}`$. Then there exist $`\theta _0\mathrm{\Lambda }_C^{}`$ and $`z^{}\stackrel{~}{Z^{}}`$ such that $`HB(z^{},\theta _0)\stackrel{~}{U}`$.
Proof : Let us recall that $`\stackrel{~}{X}/C^{}Z^{}=UV`$ where $`U=\pi \left(\stackrel{~}{U}\right)`$ $`V=\pi \left(\stackrel{~}{V}\right)`$ and $`\pi :\stackrel{~}{X}\stackrel{~}{X}/C^{}`$ is the projection .
We know that $`U`$ and $`V`$ are unbounded. Let $`x_n`$ a sequence of points in $`U`$ such that $`dist(x_n,Z^{})\mathrm{}`$. Let $`z_nZ^{}`$ such that $`dist(x_n,Z^{})=dist(x_n,z_n)`$. We consider a fundamental domain $`D\stackrel{~}{Z^{}}`$ of $`C^{}`$. There exist lifts $`\stackrel{~}{z}_nD`$ and $`\stackrel{~}{x}_n\stackrel{~}{U}`$ such that $`dist(\stackrel{~}{x}_n,Z^{})=dist(\stackrel{~}{x}_n,\stackrel{~}{z}_n)`$ tends to infinity.
By compactness we can assume that a subsequence $`\stackrel{~}{x}_{n_j}`$ converges to a point $`\theta _0\stackrel{~}{X}`$ and $`\stackrel{~}{z}_{n_j}`$ also converges to a point $`\stackrel{~}{z}\overline{D}`$. Furthermore the sequence of open balls $`B(\stackrel{~}{x}_{n_j},dist(\stackrel{~}{x}_{n_j},\stackrel{~}{z}_{n_j}))\stackrel{~}{U}`$ converges to the open horoball $`HB(\theta _0,\stackrel{~}{z})\stackrel{~}{U}`$. $`\mathrm{}`$
###### Proposition 5.13.
$`\mathrm{\Lambda }_C^{}=E\left(\mathrm{}\right)`$ and $`\mathrm{\Lambda }_C=E(\mathrm{})`$.
We first describe the idea of the proof and next state some facts we will need in order to do it.
As $`\mathrm{\Lambda }_C^{}\mathrm{\Lambda }_CE\left(\mathrm{}\right)`$ the proposition boils down to proving that $`\mathrm{\Lambda }_C^{}=E\left(\mathrm{}\right)`$. Let us assume $`\mathrm{\Lambda }_C^{}E\left(\mathrm{}\right)`$ and find a contradiction.
We will show that for any sequence $`\theta _i`$ converging to $`\theta _0`$ the geodesic starting at a point $`o`$ such that $`\left|Jac_{n1}\stackrel{~}{F}\left(o\right)\right|=1`$ and ending at $`\theta _i`$ crosses the hypersurface $`\stackrel{~}{Z^{}}`$ in a point $`z_i`$. For an appropriate choice of such a sequence $`\theta _i`$ (roughly speaking, the sequence $`\theta _i`$ is chosen to be converging to $`\theta _0`$ ”transversally to $`\mathrm{\Lambda }_C^{}`$”), the shadow (defined below) projected from $`o`$ through some geodesic ball $`B(z_i,r)`$ will not intersect $`\mathrm{\Lambda }_C^{}`$. On the other hand this shadow has to meet the limit set $`\mathrm{\Lambda }_C^{}`$ because of the shadow lemma of D. Sullivan, which leads to a contradiction.
Precisely, by Lemma (5.10) and Lemma (5.12) we know that $`\stackrel{~}{U}=\mathrm{\Lambda }_C^{}`$ and that there exists an open horoball $`HB(\theta _0,\stackrel{~}{z})\stackrel{~}{U}`$ centered at a point $`\theta \mathrm{\Lambda }_C^{}`$, and whose closure contains a point $`\stackrel{~}{z}\stackrel{~}{Z^{}}`$.
Let $`o\stackrel{~}{X}`$ and $`ET_o\stackrel{~}{X}`$ an hyperplane such that $`\stackrel{~}{F}\left(o\right)=o`$, $`\left|Jac_E\stackrel{~}{F}\left(o\right)\right|=1`$, and $`\mathrm{\Lambda }_C^{}E\left(\mathrm{}\right)`$ where $`E\left(\mathrm{}\right)`$ is the topological equator associated to $`E`$.
For each $`\theta \stackrel{~}{X}`$ we denote by $`\alpha _\theta `$ the geodesic starting from $`o`$ and such that $`\alpha _\theta \left(+\mathrm{}\right)=\theta `$.
Let $`\theta _i\stackrel{~}{X}E\left(\mathrm{}\right)=\stackrel{~}{V}\stackrel{~}{U}`$ be a sequence converging to $`\theta _0`$. By continuity, for each $`i`$ large enough, the geodesic $`\alpha _{\theta _i}`$ spends some time inside the horoball $`HB(\theta _0,\stackrel{~}{z})\stackrel{~}{U}`$ and ends up inside $`\stackrel{~}{V}`$ because $`\theta _i`$ converges to $`\theta _0`$ and $`\theta _i`$ belongs to $`\stackrel{~}{V}\stackrel{~}{U}`$.
Thus $`\alpha _{\theta _i}`$ eventually crosses $`\stackrel{~}{Z^{}}`$. Let $`z_i\alpha _{\theta _i}\stackrel{~}{Z^{}}`$. As $`\stackrel{~}{Z^{}}/C^{}`$ is compact, there is an element $`\gamma _iC^{}`$ such that $`z_i=\gamma _i\left(x_i\right)`$ where $`x_i`$ is a point in the closure $`\overline{D}`$ of a fundamental domain $`D`$ for the action of $`C^{}`$ on $`\stackrel{~}{Z^{}}`$. The points $`\gamma _i\left(x_i\right)`$ and $`\gamma _i\left(o\right)`$ stay at bounded distance because $`dist(\gamma _i\left(x_i\right),\gamma _i\left(o\right))=dist(x_i,o)dist(o,D)+diamD`$. In particular, $`lim_i\mathrm{}\gamma _i\left(o\right)=\theta _0`$.
We have proved the
###### Lemma 5.14.
Let $`\theta _i\stackrel{~}{X}E(\mathrm{})`$ be a sequence which converges to $`\theta _0`$. There exists a constant $`A`$ such that for $`i`$ large enough there exists $`z_i\stackrel{~}{Z^{}}\alpha _{\theta _i}`$ and $`\gamma _iC^{}`$ such that $`dist(z_i,\gamma _i(o))A`$ and both $`z_i`$ and $`\gamma _i(o)`$ converge to $`\theta _0`$.
Let $`x`$ and $`y`$ two points in $`\stackrel{~}{X}`$.
We define the shadow $`𝒪(x,y,R)\stackrel{~}{X}`$ of the ball $`B(y,R)`$ enlighted from the point $`x`$ by
(5.29)
$$𝒪(x,y,R)=\left\{\alpha \left(+\mathrm{}\right)\right\}$$
where $`\alpha `$ runs through the set of geodesic rays starting from $`x`$ and meeting $`B(y,R)`$.
Let $`\left\{\mu _x\right\}_x`$ be a family of Patterson measures associated to the discrete group $`C^{}`$ with critical exponent $`\delta ^{}=\delta \left(C^{}\right)`$.
The following shadow lemma is due to D.Sullivan.
###### Lemma 5.15.
, , . There exist positive constants $`C`$ and $`R`$ such that for any $`y`$ in $`\stackrel{~}{X}`$, $`\nu _y(𝒪(y,\gamma (y),R))Ce^{\delta ^{}d(y,\gamma (y))}`$
###### Corollary 5.16.
Let $`z_i`$ be defined in lemma (5.14), then we have $`𝒪(o,z_i,R+A)\mathrm{\Lambda }_C^{}\mathrm{}`$ for $`i`$ large enough.
We now prove that for a good choice of $`\theta _i`$, the shadow $`𝒪(o,z_i,R+A)`$ (with $`z_i`$ associated to $`\theta _i`$ as in lemma 5.14) never meet $`\mathrm{\Lambda }_C^{}`$ for all large $`i`$’s, ie. for any $`\theta \mathrm{\Lambda }_C^{}`$ the geodesic $`\alpha _\theta `$ does not cross $`B(z_i,R+A)`$. We have no control on the radius $`R`$ coming from the shadow lemma nor on the constant $`A`$ but we will show
###### Proposition 5.17.
There exists a sequence $`\theta _i\stackrel{~}{X}\mathrm{\Lambda }_C^{}`$ such that $`\theta _i`$ converges to $`\theta _0`$ and
$$lim_i\mathrm{}inf_{\theta \mathrm{\Lambda }_C^{}}dist(z_i,\alpha _\theta )=+\mathrm{}$$
where $`z_i=\stackrel{~}{Z}^{}\alpha _{\theta _i}`$ has been constructed in lemma (5.14).
###### Corollary 5.18.
For $`i`$ large enough, $`𝒪(o,z_i,R+A)\mathrm{\Lambda }_C^{}=\mathrm{}`$.
The corollary (5.16) and the corollary (5.18) lead to a contradiction, which ends the proof of the proposition (5.13).
The end of the paragraph is devoted to proving the proposition (5.17).
###### Lemma 5.19.
Let $`\theta _i`$ be a sequence of points in $`\stackrel{~}{X}`$ converging to $`\theta _0`$ and $`z_i`$ constructed in lemma (5.14). Assume that
$`liminf_i\mathrm{}inf_{\theta \mathrm{\Lambda }_C^{}}dist(z_i,\alpha _\theta )=C<+\mathrm{}`$ then $`lim_i\mathrm{}\frac{d(\theta _i,\mathrm{\Lambda }_C^{})}{d(\theta _i,\theta _0)}=0`$.
Proof : We first show that
(5.30)
$$\underset{i\mathrm{}}{lim}dist(z_i,\alpha _{\theta _0})=\mathrm{}$$
Recall that, for any $`z\stackrel{~}{X}`$ and $`\theta \stackrel{~}{X}`$, $`B(z,\theta )`$ equals the decreasing limit as $`t`$ tends to infinity of $`dist(z,\alpha _\theta \left(t\right))dist(o,\alpha _\theta \left(t\right))`$ where $`\alpha _\theta \left(t\right)`$ is the geodesic ray joigning $`o`$ to $`\theta `$. Therefore, as the points $`z_i\stackrel{~}{Z}`$ belongs to the complementary of the fixed horoball $`HB(\stackrel{~}{z},\theta _0)`$, we have,
(5.31)
$$dist(z_i,\alpha _{\theta _0}\left(T_i\right)T_i+B(\stackrel{~}{z},\theta _0)$$
where $`dist(z_i,\alpha _{\theta _0}\left(T_i\right))=dist(z_i,\alpha _{\theta _0})`$.
On the other hand, as $`z_i`$ tends to $`\theta _0`$, $`T_i`$ tends to infinity so (5.30) is proven.
Let $`t_i`$ be such that $`z_i=\alpha _{\theta _i}\left(t_i\right)`$. By (5.30) we have
(5.32)
$$\underset{i\mathrm{}}{lim}dist(\alpha _{\theta _i}\left(t_i\right),\alpha _{\theta _0}\left(t_i\right))=\mathrm{}.$$
Let $`u_i`$ be such that
(5.33)
$$\mathrm{dist}(\alpha _{\theta _\mathrm{i}}\left(\mathrm{u}_\mathrm{i}\right),\alpha _{\theta _0}\left(\mathrm{u}_\mathrm{i}\right))=1,$$
then in particular $`u_it_i`$ for $`i`$ large enough and by the triangle inequality we have
(5.34)
$$dist(\alpha _{\theta _i}\left(t_i\right),\alpha _{\theta _0}\left(t_i\right))2\left(t_iu_i\right)+1.$$
By (5.32) we get
(5.35)
$$\underset{i\mathrm{}}{lim}\left(t_iu_i\right)=+\mathrm{}.$$
Let us assume there exists a sequence $`\theta _i^{}\mathrm{\Lambda }_C^{}`$ and a constant $`C`$ such that
(5.36)
$$dist(z_i,\alpha _{\theta _i^{}})C<+\mathrm{}.$$
We can assume that $`C1`$.
Let $`v_i`$ be such that
(5.37)
$$dist(z_i,\alpha _{\theta _i^{}})=dist(z_i,\alpha _{\theta _i^{}}\left(v_i\right)).$$
By triangle inequality,
(5.38)
$$\left|t_iv_i\right|C$$
and,
(5.39)
$$dist(\alpha _{\theta _i^{}}\left(t_i\right),\alpha _{\theta _i}\left(t_i\right))2C.$$
On the other hand, as the curvature of $`\stackrel{~}{X}`$ is bounded above by $`1`$, a classical comparison theorem gives for any $`t[0,t_i]`$,
(5.40)
$$sinh\left(\frac{dist(\alpha _{\theta _i^{}}\left(t\right),\alpha _{\theta _i}\left(t\right))}{2}\right)sinhC.\frac{sinht}{sinht_i}.$$
Let $`s_i`$ be such that
(5.41)
$$dist(\alpha _{\theta _i^{}}\left(s_i\right),\alpha _{\theta _i}\left(s_i\right))=1.$$
There are two cases. Either $`s_it_i`$ or $`s_i<t_i`$. If $`s_i<t_i`$, we get from (5.39) and (5.40) the existence of a constant $`A`$ such that for any $`i`$,
(5.42)
$$s_it_iA,$$
and this inequality also holds when $`s_it_i`$.
¿From (5.42) we get
(5.43)
$$\frac{d(\theta _i,\theta _i^{})}{d(\theta _i,\theta _0)}=e^{s_i+u_i}e^Ae^{t_i+u_i},$$
therefore, thanks to (5.35) we obtain
(5.44)
$$lim_i\mathrm{}\frac{d(\theta _i,\theta _i^{})}{d(\theta _i,\theta _0)}=0.$$
which ends the proof of lemma (5.19). $`\mathrm{}`$
###### Lemma 5.20.
Let us assume that for every sequence $`\theta _i`$ of points in $`\stackrel{~}{X}`$ converging to $`\theta _0`$, $`lim_i\mathrm{}\frac{d(\theta _i,\mathrm{\Lambda }_C^{})}{d(\theta _i,\theta _0)}=0`$. Let $`\lambda _k\mathrm{}`$ be such that the sequence of spaces $`(\stackrel{~}{X},\lambda _kd,\theta _0)`$ converges to a space $`(S,\overline{d},0)`$ in the pointed Gromov-Hausdorff topology, then the sequence of spaces $`(\mathrm{\Lambda }_C^{},\lambda _kd,\theta _0)`$ also converges to $`(S,\overline{d},0)`$.
Proof : Let us define
$$r\left(ϵ\right)=:sup\{\frac{d(\theta ,\mathrm{\Lambda }_C^{})}{d(\theta ,\theta _0)},\theta \theta _0,d(\theta ,\theta _0)ϵ\}.$$
The assumption says that
(5.45)
$$lim_{ϵ0}r\left(ϵ\right)=0.$$
For an arbitrary metric space $`(Y,d)`$ and $`Y^{}`$ a subset of $`Y`$, let us denote $`B_{(Y,d)}(y,R)`$ the closed ball of $`(Y,d)`$ of radius $`R`$ centered at $`yY`$, and $`𝒰_ϵ^{(Y,d)}\left(Y^{}\right)`$ the $`ϵ`$-neighborhood of $`Y^{}`$ in $`(Y,d)`$. For a metric space $`(Y,d)`$ and a positive number $`\lambda `$, let us denote $`\lambda Y`$ the rescaled space $`(Y,\lambda d)`$.
By definition of the function $`r`$, we have for any $`R`$,
$$B_{\lambda _k\stackrel{~}{X}}(\theta _0,R)𝒰_{ϵ_k}^{\lambda _k\stackrel{~}{X}}B_{\lambda _k\mathrm{\Lambda }_C^{}}(\theta _0,R+ϵ_k)$$
(5.46)
$$B_{\lambda _k\stackrel{~}{X}}(\theta _0,R+2ϵ_k).$$
where $`ϵ_k=:Rr(R/\lambda _k)`$.
Let us fix $`\alpha >0`$. By definition 5.7, for any $`R>0`$, $`ϵ>0`$, there exist a map $`f:B_{\lambda _k\stackrel{~}{X}}(\theta _0,R+\alpha )S`$ such that for $`kk_0`$,
(i) $`f\left(\theta _0\right)=0`$,
for any two points $`x`$ and $`y`$ in $`B_{\lambda _k\stackrel{~}{X}}(\theta _0,R+\alpha )`$ ,
(ii) $`\left|\overline{d}(f\left(x\right),f\left(y\right))\lambda _kd(x,y)\right|ϵ`$,
and
(iii) $`B_{(S,\overline{d})}(0,R+\alpha ϵ)𝒰_ϵ^{(S,\overline{d})}f\left(B_{\lambda _k\stackrel{~}{X}}(\theta _0,R+\alpha )\right)`$.
Moreover let us prove:
(iv) $`B_{(S,\overline{d})}(0,R2ϵ)𝒰_ϵ^{(S,\delta )}f\left(B_{\lambda _k\stackrel{~}{X}}(\theta _0,R)\right)`$.
Indeed, let $`zB_{(S,\overline{d})}(0,R2ϵ)`$. By (iii), there exists $`\theta B_{\lambda _k\stackrel{~}{X}}(\theta _0,R+\alpha )`$ such that $`\overline{d}(z,f\left(\theta \right))ϵ`$. Since $`\overline{d}(z,0)R2ϵ`$, we thus deduce from triangle inequality $`\overline{d}(f\left(\theta \right),0)Rϵ`$, and therefore we get, thanks to (i) and (ii), $`\lambda _kd(\theta ,\theta _0)R`$. $`\mathrm{}`$
By (5.45), for $`ϵ`$ small enough, there exists $`k_1k_0`$ such that for any $`kk_1`$, then $`2ϵ_kϵ`$ and
$$B_{\lambda _k\stackrel{~}{X}}(\theta _0,R+2ϵ_k)B_{\lambda _k\stackrel{~}{X}}(\theta _0,R+\alpha ).$$
Therefore, by (5.46) and the above properties (i),(ii),(iii) and (iv) of the map $`f`$, and the triangle inequality we get,
$$B_{(S,\overline{d})}(0,R2ϵ)𝒰_ϵ^{(S,\overline{d})}f\left(B_{\lambda _k\stackrel{~}{X}}(\theta _0,R)\right)$$
$$𝒰_ϵ^{(S,\overline{d})}f\left(𝒰_{ϵ_k}^{\lambda _k\stackrel{~}{X}}B_{\lambda _k\mathrm{\Lambda }_C^{}}(\theta _0,R+ϵ_k)\right)$$
$$𝒰_{2ϵ+ϵ_k}^{(S,\overline{d})}f\left(B_{\lambda _k\mathrm{\Lambda }_C^{}}(\theta _0,R+ϵ_k)\right).$$
About the second inclusion above let us remark that the set $`𝒰_{ϵ_k}^{\lambda _k\stackrel{~}{X}}B_{\lambda _k\mathrm{\Lambda }_C^{}}(\theta _0,R+ϵ_k)`$ is contained in $`B_{\lambda _k\stackrel{~}{X}}(\theta _0,R+2ϵ_k)B_{\lambda \stackrel{~}{X}}(\theta _0,R+\alpha )`$, so that we can apply $`f`$ to this set.
¿From the above inclusions we obtain
$$B_{(S,\overline{d})}(0,R3ϵ)𝒰_{2ϵ+ϵ_k}^{(S,\overline{d})}f\left(B_{\lambda _k\mathrm{\Lambda }_C^{}}(\theta _0,R)\right)$$
which implies the convergence of $`(\mathrm{\Lambda }_C^{},\lambda _kd,\theta _0)`$ to $`(S,\overline{d},0)`$.
$`\mathrm{}`$
###### Corollary 5.21.
Let us assume that for every sequence $`\theta _i`$ of points in $`\stackrel{~}{X}`$ converging to $`\theta _0`$ and $`z_i`$ the sequence of points constructed in lemma (5.14), $`liminf_i\mathrm{}inf_{\theta \mathrm{\Lambda }_C^{}}dist(z_i,\alpha _\theta )<+\mathrm{}`$. Let $`\lambda _k\mathrm{}`$ be such that the sequence of spaces $`(\stackrel{~}{X},\lambda _kd,\theta _0)`$ converges to the space $`(S,\delta ,0)`$ in the pointed Gromov-Hausdorff topology, then the sequence of spaces $`(\mathrm{\Lambda }_C^{},\lambda _kd,\theta _0)`$ also converges to $`(S,\delta ,0)`$.
We will show now that there exist a sequence of points $`\theta _1^k,\theta _2^k\mathrm{\Lambda }_C^{}`$ converging to $`\theta _0`$, such that the mutual distances $`d(\theta _1^k,\theta _2^k)`$, $`d(\theta _1^k,\theta _0)`$, $`d(\theta _2^k,\theta _0)`$ is tending to zero at the same rate, and the triple $`\theta _1^k,\theta _2^k,\theta _0`$ can be uniformly separated by elements $`\gamma _kC^{}`$.
###### Lemma 5.22.
Assume that every weak tangent of $`(\stackrel{~}{X},d)`$ at $`\theta _0`$ belongs to $`WT(\mathrm{\Lambda }_C^{},d)`$, then there exist positive constants $`c`$, $`\delta `$, a sequence $`ϵ_k`$ tending to $`0`$ when $`k`$ tends to $`\mathrm{}`$, a sequence $`\gamma _kC^{}`$, a sequence of points $`\theta _1^k,\theta _2^k\mathrm{\Lambda }_C^{}`$ such that for $`i=1,2`$,
$`c^1ϵ_kd(\theta _1^k,\theta _2^k)cϵ_k`$,
$`c^1ϵ_kd(\theta _i^k,\theta _0))cϵ_k`$ and
$`d(\gamma _k\theta _1^k,\gamma _k\theta _2^k)\delta `$, $`d(\gamma _k\theta _i^k,\gamma _k\theta _0))\delta `$.
Proof : For any $`x\stackrel{~}{X}\stackrel{~}{X}`$ and $`y\stackrel{~}{X}\stackrel{~}{X}`$ let us define $`\alpha _{x,y}`$ the geodesic ray joining $`x`$ and $`y`$. Let $`o\stackrel{~}{X}`$ and $`ET_o\stackrel{~}{X}`$ be such that $`\left|Jac_E\stackrel{~}{F}\left(o\right)\right|=1`$ and $`E\left(\mathrm{}\right)`$ the equator associated to $`E`$. Let $`\gamma _kC^{}`$ be a sequence such that $`\gamma _k\left(o\right)`$ converges to the point $`\theta _0`$ where $`\theta _0\mathrm{\Lambda }_C^{}`$ is the point coming from lemma 5.12. In particular, according to that lemma, there exist a point $`\stackrel{~}{z}\stackrel{~}{Z^{}}`$ such that the hypersurface $`\stackrel{~}{Z^{}}`$ is contained in the complementary of the open horoball $`HB(\stackrel{~}{z},\theta _0)`$. We define $`D:=dist(\stackrel{~}{z},o)`$. As $`\stackrel{~}{Z^{}}`$ lies outside the open horoball $`HB(\stackrel{~}{z},\theta _0)`$, the points $`\gamma _k\left(o\right)`$ belong to the complementary of the open horoball $`HB(\alpha _{\stackrel{~}{z},\theta _0}\left(D\right),\theta _0)`$. By standard triangle comparison argument (comparison with the hyperbolic case) the angle $`Angle(\alpha _{\gamma _k\left(o\right),\theta _0},\alpha _{\gamma _k\left(o\right),o})`$ between the two geodesic rays $`\alpha _{\gamma _k\left(o\right),\theta _0}`$ and $`\alpha _{\gamma _k\left(o\right),o}`$ satisfies :
(5.47)
$$lim_k\mathrm{}Angle(\alpha _{\gamma _k\left(o\right),\theta _0},\alpha _{\gamma _k\left(o\right),o})=0.$$
By equivariance we have $`\mathrm{\Lambda }_C^{}\left(\gamma _kE\right)\left(\mathrm{}\right)`$ where $`\gamma _kET_{\gamma _k\left(o\right)}\stackrel{~}{X}`$. For any $`vT\stackrel{~}{X}`$ let $`\alpha _v`$ be the geodesic ray such that $`\dot{\alpha }_v\left(0\right)=v`$. Let us denote by $`u_k`$ the unit vector in $`\gamma _kE`$ such that $`\alpha _{u_k}\left(+\mathrm{}\right)=\theta _0`$ and let us choose some $`w_k\gamma _kE`$ such that $`<u_k,w_k>=0`$ (this is possible because $`n12`$.
We claim now that there exist $`v_k\gamma _kE`$ such that the angle between $`v_k`$ and $`w_k`$ is not too far from $`0`$ or $`\pi `$, namely
(5.48)
$$|<v_k,w_k>|\frac{1}{\left(n1\right)^{1/2}},$$
and $`\alpha _{v_k}\left(+\mathrm{}\right)\mathrm{\Lambda }_C^{}`$ or $`\alpha _{v_k}\left(\mathrm{}\right)\mathrm{\Lambda }_C^{}`$.
Let us prove this claim.
According to Proposition 5.1 and to (5.11), the restriction to $`\gamma _kE`$ of the quadratic form $`h\left(u\right)=DB(\gamma _k\left(o\right),\theta )\left(u\right)^2𝑑\mu _{\gamma _k\left(o\right)}\left(\theta \right)`$ verifies
(5.49)
$$h_{\gamma _kE}\left(u\right)=\frac{u^2}{n1}.$$
Therefore, if for all $`u\gamma _k\left(E\right)`$ such that $`\alpha _u\left(+\mathrm{}\right)=\theta \mathrm{\Lambda }_C^{}`$ we had $`|<u,w_k>|<\frac{1}{\left(n1\right)^{1/2}}`$, then one would get $`h\left(w_k\right)<\frac{1}{n1}`$, which contradicts (5.49) and proves the claim.
In particular the angle between $`u_k`$ and $`v_k`$ is not too far from $`\pi /2`$ for $`k`$ large enough, ie.
(5.50)
$$|<u_k,v_k>|\left(\frac{n2}{n1}\right)^{1/2},$$
and thanks to (5.47), we have for $`k`$ large enough
(5.51)
$$|<\dot{\alpha }_{\gamma _k\left(0\right),o}\left(0\right),v_k>|\left(\frac{n\frac{3}{2}}{n1}\right)^{1/2}.$$
Let us now assume for example that $`\theta _k=\alpha _{v_k}\left(+\mathrm{}\right)\mathrm{\Lambda }_C^{}`$. Let us show that
(5.52)
$$lim_k\mathrm{}d(\theta _0,\theta _k)=0.$$
Assume that (5.52) is not true. Then, one can assume after extracting a subsequence that $`\theta _k`$ converges to $`\theta \theta _0`$. Therefore the geodesic rays $`\alpha _{\gamma _k\left(0\right),o}`$ and $`\alpha _{v_k}`$ would converge to the geodesics $`\alpha _{\theta _0,o}`$ and $`\alpha _{\theta _0,\theta }`$ and thus the angle $`Angle(\alpha _{\gamma _k\left(0\right),o},\alpha _{v_k})`$ would converge to $`0`$. But this would contradict (5.51).
Let us now denote $`ϵ_k=:d(\theta _k,\theta _0)`$. According to (5.52), $`lim_k\mathrm{}ϵ_k=0`$. We now consider the following sequence of pointed metric space $`(\stackrel{~}{X},ϵ_k^1d,\theta _0)`$, a subsequence of which being converging to some metric space $`(S,\delta )`$, cf\[\]. For convenience we still denote by the same index $`k`$ the subsequence. By the corollary 5.21, the sequence $`(\mathrm{\Lambda }_C^{},ϵ_k^1d,\theta _0)`$ also converges to $`(S,\delta )`$. According to lemma 5.8, the space $`S`$ is homeomorphic to $`^{n1}`$. In particular there exist a sequence of points $`\theta _k^{}\mathrm{\Lambda }_C^{}`$ and a constant $`c`$ such that
(5.53)
$$c^1ϵ_kd(\theta _k,\theta _k^{})cϵ_k,$$
(5.54)
$$c^1ϵ_kd(\theta _k^{},\theta _0))cϵ_k.$$
Thus, the points $`\theta _1^k=\theta _k`$ and $`\theta _2^k=\theta _k^{}`$ satisfy the two first properties of lemma 5.22.
In order to complete the proof of lemma 5.22, we will show that the elements $`\eta _k=:\gamma _k^1`$ uniformly separate $`\theta _0`$, $`\theta _1^k`$ and $`\theta _2^k`$.
Thanks to (5.50) the angle at $`\gamma _k\left(o\right)`$ between $`\theta _1^k`$ and $`\theta _0`$ is uniformly bounded away from $`0`$ and $`\pi `$ and so does the angle at $`o`$ between $`\gamma _k^1\left(\theta _1^k\right)`$ and $`\gamma _k^1\left(\theta _0\right)`$. Therefore, as the angle is Hölder-equivalent to the distance $`d`$, cf. , there is a constant $`c`$ such that
(5.55)
$$d(\gamma _k^1\left(\theta _1^k\right),\gamma _k^1\left(\theta _0\right))c.$$
Now the cocompact group $`\mathrm{\Gamma }`$ acts uniformly quasi-conformally on $`(\stackrel{~}{X},d)`$, ( and Theorem 5.2), and so does $`C^{}\mathrm{\Gamma }`$, therefore
(5.56)
$$d(\gamma _k^1\left(\theta _1^k\right),\gamma _k^1\left(\theta _2^k\right))c,$$
and
(5.57)
$$d(\gamma _k^1\left(\theta _2^k\right),\gamma _k^1\left(\theta _0\right))c.$$
which ends the proof of lemma 5.22. $`\mathrm{}`$
Proof of Proposition 5.17 :
Let us assume that for every sequence $`\theta _i`$ of points in $`\stackrel{~}{X}`$ converging to $`\theta _0`$, $`liminf_i\mathrm{}inf_{\theta \mathrm{\Lambda }_C^{}}dist(z_i,\alpha _\theta )<+\mathrm{}`$, then by corollary 5.21 and lemma 5.22 there exist a positive constant $`c`$, a sequence $`ϵ_k`$ tending to $`0`$ when $`k`$ tends to $`\mathrm{}`$, a sequence $`\gamma _kC^{}`$, a sequence of points $`\theta _1^k,\theta _2^k\mathrm{\Lambda }_C^{}`$ such that
$`c^1ϵ_kd(\theta _1^k,\theta _2^k)cϵ_k`$,
$`c^1ϵ_kd(\theta _i^k,\theta _0))cϵ_k`$ and
$`d(\gamma _k\theta _1^k,\gamma _k\theta _2^k)\delta `$, $`d(\gamma _k\theta _i^k,\gamma _k\theta _0))\delta `$.
Applying lemma 5.9 for $`=\mathrm{\Lambda }_C^{}`$ and $`\lambda _k=ϵ_k^1`$ we conclude that $`\mathrm{\Lambda }_C^{}`$ is homeomorphic to $`\stackrel{~}{X}`$, which is impossible because $`\mathrm{\Lambda }_C^{}`$ is contained in a topological equator $`E\left(\mathrm{}\right)`$. $`\mathrm{}`$
Step 4 : $`C^{}`$ and $`C`$ are convex cocompact.
We first define convex cocompactness. For a discrete group $`C`$ of isometries acting on a Cartan Hadamard manifold of negative sectional curvature with limit set $`\mathrm{\Lambda }_C`$, one defines the geodesic hull $`𝒢\left(\mathrm{\Lambda }_C\right)`$ of $`\mathrm{\Lambda }_C`$ as the set of all geodesics both ends of whose belong to $`\mathrm{\Lambda }_C`$.
The geodesic hull of $`\mathrm{\Lambda }_C`$ is a $`C`$ invariant set. One says that $`C`$ is convex cocompact if $`𝒢\left(\mathrm{\Lambda }_C\right)/C`$ is compact.
###### Lemma 5.23.
C’ is convex cocompact.
Proof : Let us denote $`\pi :\stackrel{~}{X}\stackrel{~}{X}/C^{}`$ the projection. Assume that $`C^{}`$ is not convex cocompact. Then, there exist a sequence $`x_n𝒢\left(C^{}\right)`$ such that $`x_n`$ tends to infinity. In particular $`dist(x_n,Z^{})+\mathrm{}`$, where $`Z^{}=\stackrel{~}{Z^{}}/C^{}`$ is the compact hypersurface which separates $`\stackrel{~}{X}/C^{}`$ in two unbounded connected components. There exist lifts $`\stackrel{~}{x}_n`$ of $`x_n`$ such that
(5.58)
$$\stackrel{~}{x}_n\theta _0\mathrm{\Lambda }_C^{}$$
(5.59)
$$dist(\stackrel{~}{x}_n,\stackrel{~}{Z^{}})=dist(\stackrel{~}{x}_n,\stackrel{~}{z}_n)$$
where $`\stackrel{~}{z}_n\stackrel{~}{Z^{}}`$ is bounded. Therefore there exist $`\stackrel{~}{z}\stackrel{~}{Z^{}}`$
such that $`HB(\stackrel{~}{z},\theta _0)\stackrel{~}{U}`$, where $`\stackrel{~}{U}`$ is one of the two connected components of $`\stackrel{~}{X}\stackrel{~}{Z^{}}`$, the other being $`\stackrel{~}{V}`$.
We recall that $``$ , $`𝒩`$ are the two connected components of $`\stackrel{~}{X}\mathrm{\Lambda }_C^{}`$. We also have $`\stackrel{~}{Z^{}}=\mathrm{\Lambda }_C^{}=E\left(\mathrm{}\right)`$, and after possibly replacing $`C^{}`$ by an index two subgroup, we can assume that $`C^{}`$ preserves $`\stackrel{~}{U}`$ and $`\stackrel{~}{V}`$.
Claim : There are the two following cases.
Either one of the two boundaries $`\stackrel{~}{U}`$ or $`\stackrel{~}{V}`$ is equal to $`\mathrm{\Lambda }_C^{}`$ (in this case the other boundary is equal to $`\stackrel{~}{X}`$), or $`\stackrel{~}{U}=\overline{}`$ and $`\stackrel{~}{V}=\overline{𝒩}`$, where $`\overline{}`$ and $`\overline{𝒩}`$ are the closure of $``$ and $`𝒩`$.
Let us prove the claim. We first remark that if there exist $`\theta \stackrel{~}{U}`$, then $`\stackrel{~}{U}`$. Namely, let $`\xi `$ be any other point in $``$ and $`\alpha `$ a continuous path in $``$ joining $`\theta `$ and $`\xi `$. Since the set $`\stackrel{~}{Z^{}}\mathrm{\Lambda }_C^{}`$ is a closed subset in $`\stackrel{~}{X}\stackrel{~}{X}`$, there exist an open connected neighborhood $`W`$ of $`\alpha `$ in $`\stackrel{~}{X}\stackrel{~}{X}`$ contained in the complementary of $`\stackrel{~}{Z^{}}\mathrm{\Lambda }_C^{}`$. Therefore, as $`W\stackrel{~}{U}\mathrm{}`$, we have $`W\stackrel{~}{X}\stackrel{~}{U}`$ and $`\xi \stackrel{~}{U}`$. Let us assume that neither $`\stackrel{~}{U}`$ nor $`\stackrel{~}{V}`$ is equal to $`\mathrm{\Lambda }_C^{}`$. Then, each boundary $`\stackrel{~}{U}`$ and $`\stackrel{~}{V}`$ contains $``$ or $`𝒩`$. But on the other hand, since the set $`\stackrel{~}{Z^{}}\mathrm{\Lambda }_C^{}`$ is closed, $`\left(\stackrel{~}{U}\mathrm{\Lambda }_C^{}\right)\left(\stackrel{~}{V}\mathrm{\Lambda }_C^{}\right)=\mathrm{}`$ thus we have $`\stackrel{~}{U}=`$ and $`\stackrel{~}{V}=𝒩`$ or the other way around and the claim is proved.
Case 1 : $`\stackrel{~}{U}=\mathrm{\Lambda }_C^{}`$ and $`\stackrel{~}{V}=\stackrel{~}{X}`$ or the other way around.
In this case, we are in the situation of the step 3, which leads to a contradiction, cf. remark 5.11.
Case 2 : $`\stackrel{~}{U}=\overline{}`$ and $`\stackrel{~}{V}=\overline{𝒩}`$.
In that case, assuming $`C^{}`$ is not convex-cocompact, there exist an open horoball $`HB(\theta _0,\stackrel{~}{z})\stackrel{~}{U}`$ where $`\theta _0\mathrm{\Lambda }_C^{}`$, $`\stackrel{~}{z}\stackrel{~}{Z^{}}`$, $`\stackrel{~}{U}=\overline{}`$ and $`\stackrel{~}{V}=\overline{𝒩}`$. We will find a contradiction in a similar way as in case 1, ie. step 3. We consider a point $`o\stackrel{~}{X}`$ and an hyperplane $`ET_o\stackrel{~}{X}`$ such that $`\left|Jac_E\stackrel{~}{F}\left(o\right)\right|=1`$ and $`\mathrm{\Lambda }_C^{}=E\left(\mathrm{}\right)`$.
Let $`\theta _i𝒩`$ be a sequence which converge to $`\theta _0`$. By continuity, for $`i`$ large enough, the geodesic ray $`\alpha _{o,\theta _i}`$ spends some time in $`HB(\theta _0,\stackrel{~}{z})\stackrel{~}{U}`$ and ends up in $`\stackrel{~}{V}`$ because $`\theta _i`$ converges to $`\theta _0`$ and $`\theta _i`$ belongs to $`𝒩=\stackrel{~}{V}\mathrm{\Lambda }_C^{}`$. Therefore, $`\alpha _{o,\theta _i}`$ eventually crosses $`\stackrel{~}{Z^{}}`$. Let $`z_i`$ be some point in $`\stackrel{~}{Z^{}}\alpha _{o,\theta _i}`$.
We will prove the following Proposition, similar to the Proposition 5.18,
###### Proposition 5.24.
There exist a sequence $`\theta _i𝒩`$ such that $`\theta _i`$ converges to $`\theta _0`$ and
$$lim_i\mathrm{}inf_{\theta \mathrm{\Lambda }_C^{}}dist(z_i,\alpha _\theta )=+\mathrm{}$$
where $`z_i\stackrel{~}{Z^{}}\alpha _{\theta _i}`$.
###### Remark 5.25.
The difference between the propositions 5.24 and 5.18 is that we are looking for a sequence $`\theta _i𝒩`$ instead of $`\theta _i\stackrel{~}{X}\mathrm{\Lambda }_C^{}`$.
Assuming the Proposition 5.24 we find a contradiction in the same way as in step 3. Namely, as $`\stackrel{~}{Z^{}}/C^{}`$ is compact, the points $`z_i\stackrel{~}{Z^{}}\alpha _{\theta _i}`$ stay at bounded distance from the $`C^{}`$-orbit of a fixed point, say, $`o`$, thus there exist a constant $`A>0`$ and elements $`\gamma _iC^{}`$ such that for any $`i`$,
(5.60)
$$dist(z_i,\gamma _io)A.$$
¿From (5.60) and the shadow lemma 5.15, we obtain $`𝒪(o,z_i,R+A)\mathrm{\Lambda }_C^{}\mathrm{}`$, and on the other hand, from the proposition 5.24, we have $`𝒪(o,z_i,R+A)\mathrm{\Lambda }_C^{}=\mathrm{}`$, which gives the contradiction. It remains to prove the Proposition 5.24.
Proof of the proposition 5.24 : We argue by contradiction, like in the proof of the proposition 5.17. Let us assume that there exist a constant $`C>0`$ such that for any sequence of points $`\theta _i𝒩`$ converging to $`\theta _0`$, $`lim_i\mathrm{}inf_{\theta \mathrm{\Lambda }_C^{}}dist(z_i,\alpha _\theta )C`$, then by lemma 5.19, we have for any such sequence $`\theta _i𝒩`$
(5.61)
$$lim_i\mathrm{}\frac{d(\theta _i,\mathrm{\Lambda }_C^{})}{d(\theta _i,\theta _0)}=0.$$
The proof of the following lemma is the same as the proof of lemma 5.20.
###### Lemma 5.26.
Let us assume that for any sequence $`\theta _i𝒩`$ conveging to $`\theta _0`$, $`lim_i\mathrm{}\frac{d(\theta _i,\mathrm{\Lambda }_C^{})}{d(\theta _i,\theta _0)}=0`$. Let $`\{\lambda _k\}`$ be a sequence of positive numbers tending to $`+\mathrm{}`$ such that the sequence of spaces $`(\stackrel{~}{X},\lambda _kd,\theta _0)`$ converges to a space $`(S,\delta ,0)`$ in the pointed Gromov-Hausdorff topology, then, $`(,\lambda _kd,\theta _0)`$ also converges to $`(S,\delta ,0)`$.
Proof : Since $`\mathrm{\Lambda }_C^{}\overline{}`$, the assumption implies that $`lim_{ϵ0}r\left(ϵ\right)=0`$ where
$$r\left(ϵ\right)=sup\{\frac{d(\theta ,)}{d(\theta ,\theta _0)},\theta \theta _0,\theta 𝒩,d(\theta ,\theta _0)ϵ\}$$
and the proof goes the same way as in lemma 5.20 replacing $`\mathrm{\Lambda }_C^{}`$ by $``$. $`\mathrm{}`$
Similarly to the lemma 5.22, we have the
###### Lemma 5.27.
Let us assume that every weak tangent of $`(\stackrel{~}{X},d)`$ at $`\theta _0`$ belongs to $`WT((,d))`$. There exist positive constant $`c`$, $`\delta `$, a sequence $`ϵ_k`$ tending to $`0`$ when $`k`$ tends to $`+\mathrm{}`$, a sequence of $`\gamma _kC^{}`$, a sequence of points $`\theta _0^k=\theta _0,\theta _1^k,\theta _2^k`$ such that for $`ij\{0,1,2\}`$,
$`c^1ϵ_kd(\theta _i^k,\theta _j^k)cϵ_k`$ and
$`d(\gamma _k\theta _i^k,\gamma _k\theta _i^k)\delta .`$
We can now end the proof of the proposition 5.24. Let us assume that there exist a constant $`C>0`$ such that for every sequence $`\theta _i`$ of points in $`𝒩`$ converging to $`\theta _0`$, $`lim_i\mathrm{}inf_{\theta \mathrm{\Lambda }_C^{}}dist(\theta _i,\alpha _\theta )C`$, then by (5.61), lemma 5.26 and lemma 5.27, there exist a sequence $`ϵ_k`$ tending to $`0`$ when $`k`$ tends to $`\mathrm{}`$, a sequence $`\gamma _kC^{}`$, a sequence of points $`\theta _0^k=\theta _0,\theta _1^k,\theta _2^k`$ such that for $`ij\{0,1,2\}`$,
$`c^1ϵ_kd(\theta _i^k,\theta _j^k)cϵ_k`$ and
$`d(\gamma _k\theta _i^k,\gamma _k\theta _i^k)\delta .`$
Applying the lemma 5.9 for $`=`$ and $`\lambda _k=ϵ_k^1`$, we conclude that $``$ is homeomorphic to $`\stackrel{~}{X}`$, which is impossible because $`\stackrel{~}{X}`$ is a sphere, and $``$ is homeomorphic to an hemisphere. This ends the proof of the proposition 5.24. $`\mathrm{}`$
###### Corollary 5.28.
C is convex cocompact.
Proof : The subgroup $`C^{}`$ of $`C`$ is convex cocompact and the limit sets of $`C^{}`$ and $`C`$ coincide by step 3, therefore $`C`$ is convex cocompact. $`\mathrm{}`$.
Step 5: C preserves a copy of the $`(n1)`$-dimensional hyperbolic space $`^{n1}`$ totally geodesically embedded in $`\stackrel{~}{X}`$.
From the steps 1-4, we know that the groups $`C`$ and $`C^{}`$ are convex cocompact, and that their limit set $`\mathrm{\Lambda }_C`$ and $`\mathrm{\Lambda }_C^{}`$ are equal to a topological equator $`E\left(\mathrm{}\right)`$.
Let us consider the essential hypersurface $`Z^{}\stackrel{~}{X}/C^{}`$. We will show that there exist a minimizing current representing the class of $`Z^{}`$ in $`H_{n1}(\stackrel{~}{X}/C^{},)`$ and that this minimizing current lifts to a totally geodesic hypersurface embedded in $`\stackrel{~}{X}`$. We will then show that this totally geodesic hypersurface is eventually hyperbolic.
We work in $`\stackrel{~}{X}/C^{}`$ and consider the essential hypersurface $`Z^{}\stackrel{~}{X}/C^{}`$. We will now prove that there exist a minimal current representing the class of $`Z^{}`$ in $`H_{n1}(\stackrel{~}{X}/C^{},)`$. Let $`\left\{Z_k\right\}`$ be a minimizing sequence of currents homologous to $`Z^{}`$. The othogonal projection onto the convex core of $`\stackrel{~}{X}/C^{}`$ is distance nonincreasing and thus volume nonincreasing. Therefore we can assume that the $`Z_k`$’s are in the the convex core of $`\stackrel{~}{X}/C^{}`$, which is compact. By (5.5), the sequence $`\left\{Z_k\right\}`$ subconverges to a minimal current $`Z_{\mathrm{}}`$ in $`\stackrel{~}{X}/C^{}`$. By (8.2), $`Z_{\mathrm{}}`$ is a manifold with possible singularities of codimension greater than or equal to 8. By corollary (4.4) and minimality we get that $`\left|Jac_{n1}\stackrel{~}{F}\left(x\right)\right|=1`$ at every regular points $`xZ_{\mathrm{}}`$. We will use the fact that $`\left|Jac_{n1}\stackrel{~}{F}\left(x\right)\right|=1`$ at every regular points $`xZ_{\mathrm{}}`$ in order to prove that $`Z_{\mathrm{}}`$ is a totally geodesic hypersurface.
###### Lemma 5.29.
Let $`x`$ and $`y`$ two distinct points in $`\stackrel{~}{X}`$ and $`E_xT_x\stackrel{~}{X}`$, $`E_yT_y\stackrel{~}{X}`$ be such that $`Jac_{n1}\stackrel{~}{F^{}}(x)=Jac_{E_x}\stackrel{~}{F^{}}(x)=1`$ and $`Jac_{n1}\stackrel{~}{F^{}}(y)=Jac_{E_y}\stackrel{~}{F^{}}(y)=1`$. Then, the geodesic $`\alpha _{x,y}`$ (resp. $`\alpha _{y,x}`$) joining $`x`$ and $`y`$ (resp. $`y`$ and $`x`$) satisfies $`\dot{\alpha }_{x,y}(0)E_x`$, (resp. $`\dot{\alpha }_{y,x}(0)E_y`$). In particular, $`\alpha _{x,y}(+\mathrm{})`$ and $`\alpha _{y,x}(+\mathrm{})`$ belong to $`\mathrm{\Lambda }_C^{}`$.
Proof: Let $`S_x`$ and $`S_y`$ be the unit spheres of $`E_x`$ and $`E_y`$. For any unit tangent vector $`uT_z\stackrel{~}{X}`$ at some point $`z`$, we define $`\theta _u\stackrel{~}{X}`$ by $`\dot{\alpha }_{z,\theta _u}\left(0\right)=u`$. By step 3, $`\mathrm{\Lambda }_C^{}=E_x\left(\mathrm{}\right)=E_y\left(\mathrm{}\right)`$, therefore for every $`uS_x`$, $`\theta _u\mathrm{\Lambda }_C^{}`$ and there exist $`vE_y`$ such that $`\theta _u=\theta _v`$. As $`E_y`$ is a vector space, $`\theta _v`$ belongs to $`\mathrm{\Lambda }_C^{}`$ therefore there exist $`wE_x`$ such that $`\theta _w=\theta _v`$. The map $`f:S_xS_x`$ defined by $`f\left(u\right)=w`$ is a continuous map. The lemma then boils down to proving that there exist $`uS_x`$ such that $`f\left(u\right)=u`$ because in that case, $`x`$, $`y`$ and $`\theta _u`$ are on the same geodesic $`\alpha _{x,\theta _u}`$.
The following properties of $`f`$ are obvious.
(i) For every $`uS_x`$, $`f\left(u\right)u`$.
(ii) $`ff=Id`$.
So $`f`$ is an involution of the sphere without fixed point and for any such map, we claim that there exist $`u`$ in the sphere such that $`f\left(u\right)=u`$. In order to prove the claim, we follow a very similar argument in , theorem 1. We argue by contradiction. Let us assume that for every $`uS_x`$, $`f\left(u\right)u`$. The map $`g:S_xS_x`$ defined by $`g\left(u\right)=\frac{f\left(u\right)+u}{f\left(u\right)+u}`$, is then well defined and continuous. Let us remark that as for every $`uS_x`$, $`f\left(u\right)u`$, then $`f`$ is homotopic to the Identity, and so is $`g`$. Moreover by (ii) we clearly have $`gf=g`$, thus the map $`g`$ factorizes through $`S_x/G_f`$ where $`G_f`$ is the group generated by the involution $`f`$. By (i) $`f`$ has no fixed point thus $`S_x/G_f`$ is a manifold and the projection $`p:S_xS_x/G_f`$ is a degre 2 map. Therefore, the induced endomorphism $`g_{}`$ on $`H_{n1}(S_x,_2)`$ is trivial, which contradicts the fact that $`g`$ is homotopic to the Identity. $`\mathrm{}`$
###### Corollary 5.30.
Let $`^{n1}\stackrel{~}{X}`$ be an hypersurface with possibly non empty boundary $`^{n1}`$, such that for any $`x^{n1}`$, $`Jac_{n1}\stackrel{~}{F^{}}(x)=Jac_{E_x}\stackrel{~}{F^{}}(x)=1`$, where $`E_x`$ is the tangent space of $`^{n1}`$ at $`x`$. Let us consider $`x^{n1}`$ such that $`dist_{\stackrel{~}{X}}(x,^{n1})=r>0`$. Then, for any $`x^{}^{n1}`$ with $`dist_{\stackrel{~}{X}}(x,x^{})<r`$, the geodesic $`\alpha _{x,x^{}}`$ joining $`x`$ and $`x^{}`$ is contained in $`^{n1}`$. In particular, $`^{n1}`$ is locally convex.
Proof of the corollary: Let us fix $`\theta \mathrm{\Lambda }_C^{}`$ and consider the vector field $`B(y,\theta )`$ in $`\stackrel{~}{X}`$. Let $`x^{n1}`$. As $`Jac_{E_x}\stackrel{~}{F^{}}\left(x\right)=1`$, we have by step 3 $`\mathrm{\Lambda }_C^{}=E_x\left(\mathrm{}\right)`$. Then, for any $`x^{n1}`$, $`B(x,\theta )`$ is tangent to $`^{n1}`$, therefore the geodesic $`\alpha _{x,\theta }`$ satisfies $`\alpha _{x,\theta }\left(t\right)^{n1}`$ for all $`t[0,r)`$. Let $`x^{}^{n1}`$. By lemma 5.29, $`\dot{\alpha }_{x,x^{}}\left(0\right)E_x`$, therefore $`\alpha _{x,x^{}}=\alpha _{x,\theta }`$ and $`\alpha _{x,x^{}}\left(t\right)^{n1}`$ for all $`t[0,r)`$. $`\mathrm{}`$
We now prove that $`Z_{\mathrm{}}`$ is a totally geodesic hypersurface in $`\stackrel{~}{X}/C^{}`$. Let us recall that $`Z_{\mathrm{}}`$ is a manifold which is smooth except at a singular subset of codimension at least $`7`$. Let us consider a lift $`\stackrel{~}{Z}_{\mathrm{}}\stackrel{~}{X}`$ of $`Z_{\mathrm{}}`$ and denote $`\stackrel{~}{Z}_{\mathrm{}}^{reg}`$ (resp. $`\stackrel{~}{Z}_{\mathrm{}}^{sing}`$) the set of regular (resp.) singular points of $`\stackrel{~}{Z}_{\mathrm{}}`$.
###### Lemma 5.31.
$`\stackrel{~}{Z}_{\mathrm{}}`$ is a totally geodesic hypersurface in $`\stackrel{~}{X}`$.
Proof: Let us consider a regular point $`x\stackrel{~}{Z}_{\mathrm{}}^{reg}`$. We shall show that for every point $`x^{}\stackrel{~}{Z}_{\mathrm{}}^{reg}`$ the geodesic segment joining $`x`$ and $`x^{}`$ is contained in $`\stackrel{~}{Z}_{\mathrm{}}`$, and as the set of regular points is dense in $`\stackrel{~}{Z}_{\mathrm{}}`$ (as the complementary of a subset of codimension at least $`8`$), this will show that $`\stackrel{~}{Z}_{\mathrm{}}`$ is totally geodesic.
We claim that there exist a sequence $`y_k\stackrel{~}{Z}_{\mathrm{}}^{reg}`$ such that $`lim_k\mathrm{}y_k=x^{}`$ and the geodesic segment joining $`x`$ and $`y_k`$ is contained in $`\stackrel{~}{Z}_{\mathrm{}}`$.
The claim immediately implies that the geodesic segment joining $`x`$ and $`x^{}`$ is contained in $`\stackrel{~}{Z}_{\mathrm{}}`$.
Let us prove the claim.
For $`y\stackrel{~}{Z}_{\mathrm{}}^{reg}`$ we consider $`\alpha _{x,y}`$ the geodesic joining $`x`$ and $`y`$ and define
(5.62)
$$t_y=inf\left\{t>0,\alpha _{x,y}\left(t\right)\stackrel{~}{Z}_{\mathrm{}}\right\}$$
As $`x`$ is a regular point, by corollary 5.30, there exist $`ϵ>0`$ such that $`t_y>ϵ`$.
In order to prove the claim, we argue by contradiction. Let us assume that there exist $`r>0`$ such that for any $`yB_{\stackrel{~}{X}}(x^{},r)\stackrel{~}{Z}_{\mathrm{}}^{reg}`$, $`t_y<dist(x,y)`$. By corollary 5.30 applied to $`\stackrel{~}{Z}_{\mathrm{}}^{reg}`$, we have $`\alpha _{x,y}\left(t_y\right)\stackrel{~}{Z}_{\mathrm{}}^{sing}`$. As the set of regular points is an open subset of $`\stackrel{~}{Z}_{\mathrm{}}`$, if $`r`$ is small enough we have $`B_{\stackrel{~}{X}}(x^{},r)\stackrel{~}{Z}_{\mathrm{}}^{reg}=B_{\stackrel{~}{X}}(x^{},r)\stackrel{~}{Z}_{\mathrm{}}`$. We choose such an $`r`$ and we consider the set $`S`$ of all singular points contained in the union of all geodesic segments joining $`x`$ to a point $`yB_{\stackrel{~}{X}}(x^{},r)Z_{\mathrm{}}`$. Let us consider the map defined on $`S`$ by
$$p\left(y\right)=\alpha _{x,y}\left(ϵ\right).$$
As we already saw, for any $`yB_{\stackrel{~}{X}}(x^{},r)\stackrel{~}{Z}_{\mathrm{}}`$, we have $`t_y>ϵ`$, therefore the map $`p`$ is distance decreasing and by assumption $`p`$ is surjective onto an open subset of the sphere and $`p\left(S\right)`$ is homeomorphic to an open subset of $`^{n1}`$, therefore the Hausdorff dimension of $`S`$ is greater than or equal to $`n1`$, which contradicts the fact that the singular set has codimension at least $`8`$ in $`\stackrel{~}{Z}_{\mathrm{}}`$. $`\mathrm{}`$
The totally geodesic hypersurface $`\stackrel{~}{Z}_{\mathrm{}}\stackrel{~}{X}`$ is preserved by $`C`$, and $`\stackrel{~}{Z}_{\mathrm{}}/C`$ is of minimal volume in its homology class.
Let us prove that $`\stackrel{~}{Z}_{\mathrm{}}`$ is isometric to the hyperbolic space $`_{}^{n1}`$.
###### Lemma 5.32.
$`\stackrel{~}{Z}_{\mathrm{}}`$ is isometric to the hyperbolic space $`_{}^{n1}`$.
Proof :
As $`\stackrel{~}{Z}_{\mathrm{}}/C`$ is of minimal volume in its homology class, we have by Proposition 5.1, for all $`x\stackrel{~}{Z}_{\mathrm{}}`$, $`Jac_{E_x}\stackrel{~}{F}\left(x\right)=1`$ and $`\stackrel{~}{F}\left(x\right)=x`$, where $`E_x`$ is the tangent space of $`\stackrel{~}{Z}_{\mathrm{}}`$ at $`x`$. Moreover, we saw in the proof of proposition 5.1 that
$$H=\frac{1}{n1}Id_{D\stackrel{~}{F}\left(x\right)\left(E_x\right)}=\frac{1}{n1}Id_{E_x},$$
therefore we get from (4.11) and $`\stackrel{~}{F}\left(x\right)=x`$, that for all $`u,vT_x\stackrel{~}{Z}_{\mathrm{}}`$,
$$_{\stackrel{~}{X}}\left[DdB_{(x,\theta )}(u,v)+DB_{(x,\theta )}\left(u\right)DB_{(x,\theta )}\left(v\right)\right]𝑑\nu _x\left(\theta \right)$$
(5.63)
$$=\stackrel{~}{g}(u,v)$$
where $`\stackrel{~}{g}`$ is the metric on $`\stackrel{~}{X}`$. As $`\stackrel{~}{Z}_{\mathrm{}}`$ is totally geodesic, the relation (5.62) remains true with the Busemann function $`B^{\stackrel{~}{Z}_{\mathrm{}}}`$ of $`\stackrel{~}{Z}_{\mathrm{}}`$ instead of the Busemann function $`B`$ of $`\stackrel{~}{X}`$:
$$_{\stackrel{~}{X}}[DdB_{(x,\theta )}^{\stackrel{~}{Z}_{\mathrm{}}}(u,v)+DB_{(x,\theta )}^{\stackrel{~}{Z}_{\mathrm{}}}\left(u\right)DB_{(x,\theta )}^{\stackrel{~}{Z}_{\mathrm{}}}\left(v\right)$$
(5.64)
$$=\stackrel{~}{g}(u,v).$$
On the other hand, as $`\stackrel{~}{Z}_{\mathrm{}}`$ is totally geodesic, its sectional curvature is less than or equal to $`1`$, thus by Rauch comparison theorem, we have
(5.65)
$$DdB^{\stackrel{~}{Z}_{\mathrm{}}}(x,\theta )+DB^{\stackrel{~}{Z}_{\mathrm{}}}(x,\theta )DB^{\stackrel{~}{Z}_{\mathrm{}}}(x,\theta )\stackrel{~}{g}_{\stackrel{~}{Z}_{\mathrm{}}}$$
for all $`\theta \stackrel{~}{Z}_{\mathrm{}}=\mathrm{\Lambda }_C`$, where $`\stackrel{~}{g}_{\stackrel{~}{Z}_{\mathrm{}}}`$ is the restriction of $`\stackrel{~}{g}`$ to $`\stackrel{~}{Z}_{\mathrm{}}`$.
As the support of the measure $`\nu _x`$ is $`\stackrel{~}{Z}_{\mathrm{}}=\mathrm{\Lambda }_C`$ (by convex cocompactness of $`C`$) and the Busemann function is continuous, we get from (5.64) and (5.65) that for all $`x\stackrel{~}{Z}_{\mathrm{}}`$ and all $`\theta \stackrel{~}{Z}_{\mathrm{}}`$
$$DdB^{\stackrel{~}{Z}_{\mathrm{}}}(x,\theta )+DB^{\stackrel{~}{Z}_{\mathrm{}}}(x,\theta )DB^{\stackrel{~}{Z}_{\mathrm{}}}(x,\theta )$$
(5.66)
$$=\stackrel{~}{g}_{\stackrel{~}{Z}_{\mathrm{}}}\left(x\right).$$
and this last relation is characteristic of the hyperbolic space. $`\mathrm{}`$
Step 6: Conclusion
So far we have shown that $`C`$ preserves a totally geodesic copy of the hyperbolic space $`_{}^{n1}\stackrel{~}{X}`$ such that $`_{}^{n1}/C`$ is compact.
Our goal now is to show that $`Y=:_{}^{n1}/C`$ injects diffeomorphically in $`X=\stackrel{~}{X}/\mathrm{\Gamma }`$ and separates $`X`$ in two connected components $`R`$ and $`S`$ such that $`A=\pi _1\left(R\right)`$ and $`B=\pi _1\left(S\right)`$.
In order to do this, we will consider the $`\mathrm{\Gamma }`$ orbit of $`_{}^{n1}`$ in $`\stackrel{~}{X}`$ and the two connected components $`U`$ and $`V`$ of $`\stackrel{~}{X}\mathrm{\Gamma }_{}^{n1}`$ which are adjacent to $`_{}^{n1}`$. The stabilizers $`\overline{A}`$, $`\overline{B}`$ and $`\overline{C}`$ of $`U`$, $`V`$ and $`_{}^{n1}`$ contain respectively $`A`$, $`B`$ and $`C`$ and the hypersurface $`_{}^{n1}/C`$ injects in $`X=\stackrel{~}{X}/\mathrm{\Gamma }`$ and separates $`X`$ in two connected components $`R`$ and $`S`$ such that $`\pi _1\left(R\right)=\overline{A}`$ and $`\pi _1\left(S\right)=\overline{B}`$. We then show that $`\overline{C}=C`$, $`\overline{A}=A`$ and $`\overline{B}=B`$.
Let $`\overline{C}`$ be the stabilizer of $`^{n1}`$, namely $`\overline{C}=\left\{\gamma \mathrm{\Gamma },\gamma ^{n1}=^{n1}\right\}`$. We have $`C\overline{C}`$ and as $`^{n1}/C`$ is compact, so is $`^{n1}/\overline{C}`$ and thus \[$`\overline{C}:C`$\] $`<\mathrm{}`$.
Let $`p:\stackrel{~}{X}/CX=\stackrel{~}{X}/\mathrm{\Gamma }`$ and $`\overline{p}:\stackrel{~}{X}/\overline{C}X=\stackrel{~}{X}/\mathrm{\Gamma }`$ the natural projections. We now show that the restriction of $`p`$ to $`^{n1}/C`$ is an embedding, thus $`Y:=p\left(^{n1}/C\right)`$ is a compact totally geodesic hypersurface of $`X`$.
In the section 2, we constructed a $`C`$-invariant hypersurface $`\stackrel{~}{Z}\stackrel{~}{X}`$ such that $`Z=\stackrel{~}{Z}/C\stackrel{~}{X}/C`$ is compact. The hypersurface is defined as $`\stackrel{~}{Z}=\stackrel{~}{f}^1\left(t_0\right)`$ where $`\stackrel{~}{f}:\stackrel{~}{X}T`$ is an equivariant map onto the Bass-Serre tree associated to the amalgamation $`A_CB`$ and $`t_0`$ belongs to that edge of $`T`$ which is fixed by $`C`$.
Let us first show two lemmas.
###### Lemma 5.33.
The restriction of $`p`$ to $`\stackrel{~}{Z}/C`$ is an embedding into $`X=\stackrel{~}{X}/\mathrm{\Gamma }`$.
Proof : Let $`\gamma \mathrm{\Gamma }`$, $`z`$, $`z^{}`$ in $`\stackrel{~}{Z}`$ such that $`z^{}=\gamma z`$. By equivariance,
$$\stackrel{~}{f}\left(\gamma z\right)=\gamma \stackrel{~}{f}\left(z\right)=\gamma t_0=\stackrel{~}{f}\left(z^{}\right)=t_0,$$
thus $`\gamma C`$.$`\mathrm{}`$
###### Lemma 5.34.
The restriction of $`\overline{p}`$ to $`^{n1}/\overline{C}`$ is an embedding into $`X=\stackrel{~}{X}/\mathrm{\Gamma }`$.
Proof : Let us assume that there is a $`\gamma \mathrm{\Gamma }\overline{C}`$ such that $`\gamma ^{n1}^{n1}\mathrm{}`$ and choose an $`x\gamma ^{n1}^{n1}`$. As $`\gamma \overline{C}`$, there exist $`uT_x\gamma ^{n1}T_x^{n1}`$. We consider $`c_u`$ the geodesic ray such that $`\dot{c}_u\left(0\right)=u`$. We know that $`\stackrel{~}{Z}`$ is contained in an $`ϵ`$-neighbourhood $`𝒰_ϵ^{n1}`$ of $`^{n1}`$. The $`ϵ`$-neighbourhood $`𝒰_ϵ^{n1}`$ of $`^{n1}`$ separates $`\stackrel{~}{X}`$ in two connected components $`U`$ and $`V`$ and for $`t>0`$ large enough, we have, say, $`c_u\left(t\right)U`$ and $`c_u\left(t\right)V`$.
Let $`\stackrel{~}{Z^{}}`$ be the connected component of $`\stackrel{~}{Z}`$ that we constructed at the end of section 2, whose stabilizer (or an index two subgroup of it) $`C^{}`$ is such that $`\stackrel{~}{Z^{}}/C^{}`$ separates $`\stackrel{~}{X}/C^{}`$ in two unbounded connected components $`U^{}/C^{}`$ and $`V^{}/C^{}`$ where $`U^{}`$ and $`V^{}`$ are the two connected components of $`\stackrel{~}{X}\stackrel{~}{Z^{}}`$.
We claim that $`UU^{}`$ and $`VV^{}`$ or the other way around. Indeed if not, $`U`$ and $`V`$ would be both contained in, say, $`U^{}`$. But in that case, $`V^{}`$ would be contained in $`𝒰_ϵH^{n1}`$ and therefore $`V^{}/C^{}`$ would be bounded, which is a contradiction.
As $`\gamma \stackrel{~}{Z}`$ lies in the $`ϵ`$ neighborhood of $`\gamma ^{n1}`$, there exist sequences $`z_k`$, $`z_k^{}`$ in $`\gamma \stackrel{~}{Z}`$ such that $`dist(z_k,c_u\left(k\right))ϵ`$ and $`dist(z_k^{},c_u\left(k\right))ϵ`$. By proposition 5.13 and lemma 5.23, $`C^{}`$ also acts cocompactly on $`^{n1}`$, thus $`C^{}`$ is of finite index in $`C`$, and therefore there are finitely many connected components of $`\stackrel{~}{Z}`$ and the same holds for $`\gamma \stackrel{~}{Z}`$. We thus can assume that the $`z_k`$’s and $`z_k^{}`$’s belong to a single connected component of $`\gamma \stackrel{~}{Z}`$. Let us consider a continuous path $`\alpha \gamma \stackrel{~}{Z}`$ joining $`z_k`$ and $`z_k^{}`$.
By construction the distance between $`c_u\left(k\right)`$ \[resp. $`c_u\left(k\right)`$\] and $`^{n1}`$ tends to infinity and thus, for $`k`$ large enough, $`z_kU`$ and $`z_k^{}V`$ or the other way around. By the claim, we then have $`z_kU^{}`$ and $`z_k^{}V^{}`$, therefore the path $`\alpha `$ has to cross $`\stackrel{~}{Z^{}}`$ which contradicts the lemma (5.29) and ends the proof of the lemma 5.30. $`\mathrm{}`$
As we already saw, $`\stackrel{~}{Z}`$ has finitely many connected components, and so does $`\stackrel{~}{X}\stackrel{~}{Z}`$. Let us write $`\left\{W_j\right\}_{j=1,..,m}`$ the connected components of $`\stackrel{~}{X}\stackrel{~}{Z}`$. As $`C`$ acts cocompactly on $`\stackrel{~}{Z}`$ and $`^{n1}`$ there exist $`ϵ>0`$ such that $`^{n1}𝒰_ϵ\stackrel{~}{Z}`$ and $`\stackrel{~}{Z}𝒰_ϵ^{n1}`$. Moreover $`𝒰_ϵ^{n1}`$ separates $`\stackrel{~}{X}`$ in two connected components $`U`$ and $`V`$.
###### Lemma 5.35.
Let us consider $`ϵ`$ such that $`\stackrel{~}{Z}𝒰_ϵ^{n1}`$ and $`U`$ and $`V`$ the two connected components of $`\stackrel{~}{X}𝒰_ϵ^{n1}`$. There are two distinct connected components $`W_1`$ and $`W_2`$ of $`\stackrel{~}{X}\stackrel{~}{Z}`$ such that $`UW_1`$ and $`VW_2`$. Moreover, $`\stackrel{~}{f}(W_1)\stackrel{~}{T}_1`$ and $`\stackrel{~}{f}(W_2)\stackrel{~}{T}_2`$, where $`\stackrel{~}{T}_1`$ and $`\stackrel{~}{T}_2`$ are the two connected components of $`\stackrel{~}{T}\{t_0\}`$.
Proof : We argue by contradiction. Let us assume that $`U`$ and $`V`$ are contained in the same connected component $`W_1`$ of $`\stackrel{~}{X}\stackrel{~}{Z}`$. Then, all other components $`W_j`$, $`j1`$, satisfy $`W_j𝒰_ϵ^{n1}𝒰_{2ϵ}\stackrel{~}{Z}`$. Therefore, as $`C`$ acts cocompactly on $`𝒰_{2ϵ}\stackrel{~}{Z}`$, there exist a constant $`D`$ such that for any $`j1`$, $`max_{wW_j}dist_{\stackrel{~}{T}}(\stackrel{~}{f}\left(w\right),t_0)D`$. Thus, $`\stackrel{~}{f}\left(W_1\right)`$ is contained in one connected component of $`\stackrel{~}{T}\left\{t_0\right\}`$ and $`\stackrel{~}{f}\left(_{j1}W_j\right)`$, contained in the ball $`B_{\stackrel{~}{T}}(t_0,D)`$ of $`\stackrel{~}{T}`$ of radius $`D`$ centered at $`t_0`$, is bounded. This is clearly impossible because $`\stackrel{~}{T}\left\{t_0\right\}`$ has two unbounded connected components and $`\stackrel{~}{f}`$ is onto. $`\mathrm{}`$
Let us denote $`𝒜=A^{n1}`$ the $`A`$-orbit of the $`C`$-invariant totally geodesic copy of the real hyperbolic space $`^{n1}`$, and $`\overline{A}`$ the stabilizer of $`𝒜`$, ie. $`\overline{A}=\left\{\gamma \mathrm{\Gamma },\gamma 𝒜=𝒜\right\}`$. We define in a similar way $`=B^{n1}`$ and $`\overline{B}=\left\{\gamma \mathrm{\Gamma },\gamma =\right\}`$.
Let us recall that $`\overline{C}`$ is the stabilizer of $`^{n1}`$ in $`\mathrm{\Gamma }`$. We now prove the following
###### Lemma 5.36.
We have $`\overline{A}=A\overline{C}`$ and $`\overline{B}=B\overline{C}`$. Moreover, $`\overline{A}`$ and $`\overline{B}`$ are charactrized by $`\overline{A}=\{\gamma \mathrm{\Gamma },\gamma ^{n1}𝒜\}`$ and $`\overline{B}=\{\gamma \mathrm{\Gamma },\gamma ^{n1}\}`$.
Proof : Let $`\gamma ^{}\overline{A}`$, then $`\gamma ^{}^{n1}𝒜`$ and thus there exist $`\gamma A`$ such that $`\gamma ^{}^{n1}=\gamma ^{n1}`$, therefore $`\gamma ^1\gamma ^{}\overline{C}`$, which proves the first part of the lemma.
Let us prove the second part of the lemma.
Let $`\gamma ^{}\mathrm{\Gamma }`$ be such that $`\gamma ^{}^{n1}𝒜`$. Then there exist $`\gamma A`$ such that $`\gamma ^{}^{n1}=\gamma ^{n1}`$, thus $`\gamma ^1\gamma ^{}\overline{C}`$ and therefore $`\gamma ^{}\gamma \overline{C}A\overline{C}=\overline{A}`$. This proves one inclusion, the other inclusion being obvious. $`\mathrm{}`$
For each $`\gamma \mathrm{\Gamma }`$, $`\gamma ^{n1}`$ separates $`\stackrel{~}{X}`$ in two connected components $`U_\gamma `$ and $`V_\gamma `$.
Let us now prove the following lemma.
###### Lemma 5.37.
(i) Let $`\gamma A`$, \[resp. $`\gamma B`$\]. Then, we have $`𝒜\left\{\gamma ^{n1}\right\}U_\gamma `$ or $`𝒜\left\{\gamma ^{n1}\right\}V_\gamma `$, \[resp. $`\left\{\gamma ^{n1}\right\}U_\gamma `$ or $`\left\{\gamma ^{n1}\right\}V_\gamma `$.\]
(ii) Let $`\gamma `$ be an element of $`\mathrm{\Gamma }\overline{A}`$, \[resp. $`\mathrm{\Gamma }\overline{B}`$\]. Then $`𝒜U_\gamma `$ or $`𝒜V_\gamma `$, \[resp. $`U_\gamma `$ or $`V_\gamma `$\].
Proof : (i) We argue by contradiction. Let us consider $`\gamma ^{n1}`$, $`\gamma ^{}^{n1}`$ and $`\gamma ^{\prime \prime }^{n1}`$ three distinct elements in $`𝒜`$ such that $`\gamma ^{}^{n1}U_\gamma `$ and $`\gamma ^{\prime \prime }^{n1}V_\gamma `$. By equivariance we can assume $`\gamma `$ is the identity. Let us recall that $`U`$ and $`V`$ are the two connected components of $`\stackrel{~}{X}𝒰_ϵ^{n1}`$.
We then have $`\gamma ^{}^{n1}U\mathrm{}`$ and $`\gamma ^{\prime \prime }^{n1}V\mathrm{}`$, which implies $`\gamma ^{}\stackrel{~}{Z}U\mathrm{}`$ and $`\gamma ^{\prime \prime }\stackrel{~}{Z}V\mathrm{}`$.
By lemma (5.31), $`UW_1`$ and $`VW_2`$ where $`W_1`$ and $`W_2`$ are two connected components of $`\stackrel{~}{X}\stackrel{~}{Z}`$ and $`\stackrel{~}{f}\left(U\right)\stackrel{~}{T}_1`$ and $`\stackrel{~}{f}\left(V\right)\stackrel{~}{T}_2`$, therefore, $`\stackrel{~}{f}\left(U\right)`$ contains $`\gamma ^{}t_0\stackrel{~}{T}_1`$ and $`\stackrel{~}{f}\left(V\right)`$ contains $`\gamma ^{\prime \prime }t_0\stackrel{~}{T}_2`$. This is impossible because for all elements $`\gamma ^{}`$ and $`\gamma ^{\prime \prime }`$ in $`A`$, $`\gamma ^{}t_0`$ and $`\gamma ^{\prime \prime }t_0`$ belong to the same connected component of $`\stackrel{~}{T}\left\{t_0\right\}`$.
(ii) Let us consider $`\gamma \mathrm{\Gamma }\overline{A}`$. We argue by contradiction. Let us assume there exist $`\gamma `$, $`\gamma ^{}`$ in $`A`$ such that
$$\gamma ^{}^{n1}U_\gamma $$
(5.67)
$$\gamma ^{\prime \prime }^{n1}V_\gamma $$
Let $`ϵ>0`$ such that $`\stackrel{~}{Z}𝒰_ϵ^{n1}`$ and $`U`$ and $`V`$ the connected component of $`\stackrel{~}{X}𝒰_ϵ^{n1}`$. By lemma (5.31) we have $`\stackrel{~}{f}\left(\gamma U\right)\gamma \stackrel{~}{T}_1`$ and $`\stackrel{~}{f}\left(\gamma V\right)\gamma \stackrel{~}{T}_2`$, where $`\gamma \stackrel{~}{T}_1`$ and $`\gamma \stackrel{~}{T}_2`$ are the two connected components of $`\stackrel{~}{T}\left\{\gamma t_0\right\}`$. By assumption (5.62), we have
$`\gamma ^{}^{n1}\gamma U\mathrm{}`$ and $`\gamma ^{\prime \prime }^{n1}\gamma V\mathrm{}`$, which implies $`\gamma ^{}\stackrel{~}{Z}\gamma U\mathrm{}`$ and $`\gamma ^{\prime \prime }\stackrel{~}{Z}\gamma V\mathrm{}`$, therefore $`\gamma ^{}t_0\gamma \stackrel{~}{T}_1`$ and $`\gamma ^{\prime \prime }t_0\gamma \stackrel{~}{T}_2`$, which is impossible because in the tree $`\stackrel{~}{T}`$, the points $`\gamma ^{}t_0`$ and $`\gamma ^{\prime \prime }t_0`$ belong to two adjacent edges.
$`\mathrm{}`$
By lemma (5.33) (i), for every $`\gamma `$ in $`\overline{A}`$, \[resp. $`\overline{B}`$ \], we can define $`U_\gamma `$ as the connected component of $`\stackrel{~}{X}^{n1}`$ which contains all $`\gamma ^{}^{n1}`$ for all $`\gamma ^{}`$ in $`\overline{A}`$, \[ resp. $`\overline{B}`$ \], and $`\gamma ^{}^{n1}\gamma ^{n1}`$.
Let us define
(5.68)
$$U_A:=_{\gamma A}U_\gamma .$$
By definition, $`U_A`$ \[resp. $`U_B`$ \]is a convex set in $`\stackrel{~}{X}`$ whose boundary is the collection $`𝒜`$, \[resp. $``$ \], of $`\gamma ^{n1}`$, $`\gamma `$ in $`\overline{A}`$ \[resp. $`\overline{B}`$ \] and by lemma (5.33) (ii), $`U_A`$ and $`U_B`$ are two disjoint connected components of $`\stackrel{~}{X}\mathrm{\Gamma }^{n1}`$.
In fact, $`U_A`$ \[resp. $`U_B`$\], is the convex hull of $`𝒜`$, \[resp. $``$\], and $`\overline{A}`$, \[resp. $`\overline{B}`$\], is the stabilizer of $`U_A`$, \[resp. $`U_B`$\].
###### Lemma 5.38.
The closures of $`U_A`$ and $`U_B`$ intersect along $`^{n1}`$ and $`\overline{A}\overline{B}=\overline{C}`$. Moreover, no element of $`\mathrm{\Gamma }`$ sends $`U_A`$ on $`U_B`$ nor the other way around.
Proof : The convex set $`U_A`$ is the intersection of open half spaces $`U_\gamma `$, $`\gamma \overline{A}`$, and is delimited by the disjoint union of hyperplanes $`\gamma ^{n1}`$, for some $`\gamma \overline{A}`$. The same is true for $`U_B`$ and as $`U_AU_B=\mathrm{}`$, the closures of $`U_A`$ and $`U_B`$ can intersect only along one of the connected components of their boundaries, thus along $`^{n1}`$ which is obviously in both closures. This proves the first part of the lemma, let us prove the second part. By lemma 5.32, $`\overline{C}\overline{A}\overline{B}`$. Conversely, let us take $`\gamma \overline{A}\overline{B}`$, then $`\gamma `$ preserves the closures of $`U_A`$ and $`U_B`$, thus it preserves their intersection $`^{n1}`$, and therefore $`\gamma \overline{C}`$.
Let us prove the last part of the lemma. Let $`\gamma `$ be an element such that $`\gamma U_A=U_B`$. As $`^{n1}`$ is one component of the boundary $``$ of $`U_B`$, there exist one component $`\gamma ^{}^{n1}𝒜`$, $`\gamma ^{}`$ being in $`\overline{A}`$, such that $`\gamma \left(\gamma ^{}^{n1}\right)=^{n1}`$. Therefore, $`\gamma \gamma ^{}\overline{C}`$, thus $`\gamma \overline{A}`$. The same argument yields $`\gamma ^1\overline{B}`$, so $`\gamma \overline{A}\overline{B}=\overline{C}`$ and $`\gamma `$ preserves $`U_A`$ and $`U_B`$, which contradicts our choice of $`\gamma `$.
$`\mathrm{}`$
The $`\mathrm{\Gamma }`$-orbit of the closure of $`U_AU_B`$ covers $`\stackrel{~}{X}`$. Let us construct a tree $`\overline{T}`$ embedded in $`\stackrel{~}{X}`$ in the following way: the set of vertices is the set of the connected components of $`\stackrel{~}{X}\mathrm{\Gamma }^{n1}`$ and two vertices are joined by an edge if the boundaries of their corresponding connected components intersect non trivially in $`\stackrel{~}{X}`$. By construction $`\mathrm{\Gamma }`$ acts on $`\overline{T}`$, the stabilizers of the vertices $`a`$ and $`b`$ corresponding to $`U_A`$ and $`U_B`$ are $`\overline{A}`$ and $`\overline{B}`$, the stabilizer of the edge between $`a`$ and $`b`$ is $`\overline{C}`$ and a fundamental domain for this action is the segment joining $`a`$ and $`b`$. By , I, 4, Theorem 6, the group $`\mathrm{\Gamma }`$ is the amalgamated product of $`\overline{A}`$ and $`\overline{B}`$ over $`\overline{C}`$.
We now claim that $`\overline{A}=A`$, $`\overline{B}=B`$ and $`\overline{C}=C`$.
As $`A`$, $`B`$ and $`C`$ are subgroups of $`\overline{A}`$, $`\overline{B}`$, and $`\overline{C}`$, the corresponding Mayer-Vietoris sequences of $`A_CB`$ and $`\overline{A}_{\overline{C}}\overline{B}`$ are related by the following commutative diagram
$$\begin{array}{cccccc}& H_n(A,)H_n(B,)& & H_n(\mathrm{\Gamma },)& & H_{n1}(C,)\\ & & & & & \\ & H_n(\overline{A},)H_n(\overline{B},)& & H_n(\mathrm{\Gamma },)& & H_{n1}(\overline{C},).\end{array}$$
We know that the index $`\left[\overline{C}:C\right]`$ is finite, and by the lemma 5.32, the indices of $`A`$, $`B`$, and $`C`$ in $`\overline{A}`$, $`\overline{B}`$ and $`\overline{C}`$ are finite and equal. On the other hand, the indices $`\left[\mathrm{\Gamma }:A\right]`$ and $`\left[\mathrm{\Gamma }:B\right]`$ are infinite by assumption, thus the previous diagram becomes
$$\begin{array}{cccccc}& 0& & H_n(\mathrm{\Gamma },)& & H_{n1}(C,)\\ & & & & & \\ & 0& & H_n(\mathrm{\Gamma },)& & H_{n1}(\overline{C},).\end{array}$$
Moreover, the map $`H_n(\mathrm{\Gamma },)H_{n1}(\overline{C},)`$ is bijective. Namely, the injectivity comes from the above diagram and the surjectivity from the fact that that the hypersurface $`^{n1}/\overline{C}`$ bounds in $`^n/A`$ and $`^n/B`$ so that the map $`H_{n1}(C,)H_{n1}(A,)H_{n1}(B,)`$ is trivial. Therefore the index $`[\overline{C}:C]=1`$ and we get $`\overline{A}=A`$, $`\overline{B}=B`$ and $`\overline{C}=C`$. $`\mathrm{}`$
## 6. Proof of the theorems 1.5 and 1.6.
The proof of theorem 1.5 is exactly the same as the proof of theorem 1.2. The actions of $`\mathrm{\Gamma }`$ on $`\stackrel{~}{(}X)`$ and $`T`$ give rise to a continuous $`\mathrm{\Gamma }`$-equivariant map $`\stackrel{~}{f}:\stackrel{~}{X}T`$. Like in section 2, we build an hypersurface $`\stackrel{~}{f}^1\left(t_0\right)`$ where $`t_0`$ is a regular value of $`\stackrel{~}{f}`$ belonging the interior of an edge. As the edge separates the tree in two unbounded components, the section 2 applies and we get a subgroup $`C^{}`$ of $`C`$, and an hypersurface $`\stackrel{~}{Z}^{}\stackrel{~}{X}/C^{}`$ which is essential. Now, if the action of $`\mathrm{\Gamma }`$ is minimal, every edge separates $`T`$ in two unbounded components. $`\mathrm{}`$
## 7. Appendix
The goal of this section is to give a proof of lemma 5.9. This lemma is contained in lemma 2.1, 5.1 and 5.2 of , but our situation being not exactly the same, we reproduce it down here for sake of completness.
Let us restate the lemma 5.9.
###### Lemma 7.1.
Let $`\stackrel{~}{X}`$ be a closed $`C^{}`$-invariant subset and $`\theta _0`$. We assume that there exist a sequence of positive real numbers $`\lambda _k\mathrm{}`$ such that the sequence of pointed metric spaces $`(,\lambda _kd,\theta _0)`$ converges in the pointed Gromov-Hausdorff topology to $`(S,\delta ,0)`$ where $`(S,\delta ,0)`$ is a weak tangent of $`(\stackrel{~}{X},d)`$. We also assume that there exist positive constants $`C`$ and $`\delta `$, a sequence of points $`\theta _0^k=\theta _0,\theta _1^k,\theta _2^k`$ and a sequence of elements $`\gamma _kC^{}`$ such that $`C^1\lambda _k^1d(\theta _i^k,\theta _j^k)C\lambda _k^1`$ and $`d(\gamma _k\theta _i^k,\gamma _k\theta _j^k)\delta `$ for all $`0ij2`$. Then, $``$ is homeomorphic to the one point compactification $`\widehat{S}`$ of $`S`$. In particular $``$ is homeomorphic to $`\stackrel{~}{X}`$.
We first give a definition of pointed Hausdorff-Gromov convergence which is equivalent to the definition 5.7. We follow , paragraph 4.
A sequence of metric spaces $`(Z_k,d_k,z_k)`$ converges to the metric space $`(S,\delta ,0)`$ if for every $`R>0`$, and every $`ϵ>0`$, there exist an integer $`N`$, a subset $`DB_S(0,R)`$, subsets $`D_kB_{Z_k}(z_k,R)`$ and bijections $`f_k:D_kD`$ such that for $`kN`$,
(i) $`f_k\left(z_k\right)=0`$,
(ii) the set $`D`$ is $`ϵ`$-dense in $`B_S(0,R)`$, and the sets $`D_k`$
are $`ϵ`$-dense in $`B_{Z_k}(z_k,R)`$,
(iii) $`\left|d_{Z_k}(x,y)d_Z(f_k\left(x\right),f_k\left(y\right))\right|<ϵ`$,
where $`x`$, $`y`$ belong to $`D_k`$.
Let us describe now the lemmas 2.1 and 5.1 following .
For a metric space $`(Z,d)`$ the cross ratio of four points $`\left\{z_i\right\}`$, $`i=1,\mathrm{}4`$, is the quantity
(7.1)
$$[z_1,z_2,z_3,z_4]:=\frac{d(z_1,z_3)d(z_2,z_4)}{d(z_1,z_4)d(z_2,z_3)}$$
Given two metric spaces $`X`$ and $`Y`$, an homeomorphism $`\eta :[0,\mathrm{})[0,\mathrm{})`$, and an injective map $`f:XY`$, we say that $`f`$ is an $`\eta `$-quasi-Möbius map if for any four points $`\left\{x_i\right\}`$, $`i=1,..,4`$, in $`X`$, we have
(7.2)
$$[f\left(x_1\right),f\left(x_2\right),f\left(x_3\right),f\left(x_4\right)]\eta \left([x_1,x_2,x_3,x_4]\right).$$
For example, any discrete cocompact group of isometries of $`\stackrel{~}{X}`$, where $`\stackrel{~}{X}`$ is a Cartan-Hadamard manifold with sectional curvature $`K1`$, is acting on the ideal boundary $`(\stackrel{~}{X},d)`$ endowed with the Gromov distance by $`\eta `$-quasi-Möbius transformations for some $`\eta `$.
###### Lemma 7.2 (, Lemma 2.1 ).
Let $`(X,d_X)`$ and $`(Y,d_Y)`$ be two compact metric spaces, and for any integer $`k`$, $`g_k:\stackrel{~}{D}_kY`$ an $`\eta `$-quasi-Möbius map defined on a subset $`\stackrel{~}{D}_k`$ of $`X`$. We assume that the Hausdorff distance between $`\stackrel{~}{D}_k`$ and $`X`$ satisfies
$$lim_k\mathrm{}dist_H(\stackrel{~}{D}_k,X)=0$$
and that for any integer $`k`$, there exist points $`(x_1^k,x_2^k,x_3^k)`$ in $`D_k`$ and $`(y_1^k,y_2^k,y_3^k)`$ in $`Y`$, such that $`g_k(x_i^k)=y_i^k`$ for $`i\{1,2,3\}`$, $`d_X(x_i^k,x_j^k)\delta `$ and $`d_Y(y_i^k,y_j^k)\delta `$ for $`i,j\{1,2,3\},ij`$, where $`\delta `$ is independant of $`k`$. Then a subsequence of $`g_k`$ converges uniformly to a quasi-Möbius map $`f:XY`$, ie. $`lim_{k_j\mathrm{}}dist_H(g_{k_j},f|_{\stackrel{~}{D}_{k_j}})=0`$. If in addition, we suppose that
$$lim_k\mathrm{}dist_H(g_k\left(\stackrel{~}{D}_k\right),Y)=0,$$
then the sequence $`\{g_{k_j}\}`$ converges uniformly to a quasi-Möbius homeomorphism $`f:XY`$.
Before stating the second lemma, let us define a metric space $`Z`$ to be uniformly perfect if there exist a constant $`\lambda 1`$ such that for every $`zZ`$ and $`0<R<diamZ`$, we have $`\overline{B}(z,R)B(z,\frac{R}{\lambda })\mathrm{}`$.
###### Lemma 7.3 (, lemma 5.1 ).
Let $`Z`$ be a compact uniformly perfect metric space and $`G`$ an $`\eta `$-quasi-Möbius action on $`Z`$. Suppose that for each integer $`k`$ we are given a set $`D_k`$ in a ball $`B_k=B(z,R_k)Z`$ that is $`(ϵ_kR_k)`$-dense in $`B_k`$, where $`ϵ_k>0`$, distinct points $`x_1^k,x_2^k,x_3^kB(z,\lambda _kR_k)`$, where $`\lambda _k>0`$, with
$$d_Z(x_i^k,x_j^k)\delta _kR_k$$
for $`i,j\{1,2,3\},ij`$, where $`\delta _k>0`$, and groups elements $`\gamma _kG`$ such that for $`y_i^k:=\gamma _k(x_i^k)`$ we have,
$$d_Z(y_i^k,y_j^k)\delta ^{}$$
for $`i,j\{1,2,3\},ij`$, where $`\delta ^{}`$ is independant of $`k`$. Let $`D_k^{}=\gamma _k(D_k)`$, and suppose that $`\lambda _k0`$ when $`k\mathrm{}`$, and the sequence $`\frac{ϵ_k}{\delta _k^2}`$ is bounded. Then $`lim_k\mathrm{}dist_H(D_k^{},Z)=0`$.
Let us go back to the proof of lemma 6.1. By definition of convergence, there exist a subsequence of $`\left\{\lambda _k\right\}`$, which we still denote by $`\left\{\lambda _k\right\}`$, subsets $`\stackrel{~}{D}_kB_S(0,k)`$, $`D_kB_{\lambda _k}(\theta _0,k)`$, where $`\stackrel{~}{D}_k`$ and $`D_k`$ are minimal $`1/k`$-dense subsets of $`B_S(0,k)`$ and $`B_{(,\lambda _kd)}(\theta _0,k)`$, and bijections $`f_k:\stackrel{~}{D}_kD_k`$ such that for all $`x,y\stackrel{~}{D}_k`$,
(7.3)
$$\frac{1}{2}\delta (x,y)\lambda _kd(f_k\left(x\right),f_k\left(y\right))2\delta (x,y),$$
cf. , (5.4).
We can suppose that the points $`\theta _0^k:=\theta _0`$, $`\theta _1^k`$, and $`\theta _2^k`$ in lemma 6.1 belong to the set $`D_k`$. By assumption there exist elements $`\gamma _kC^{}`$ and a constant $`\delta `$ such that
(7.4)
$$d(\gamma _k\theta _i^k,\gamma _k\theta _j^k)\delta $$
for all $`i,j\{0,1,2\}`$.
The lemma 6.1 is a direct consequence of the lemma 6.2 applied to $`(X,d_X)=(\widehat{S},\widehat{\delta })`$ and $`(Y,d_Y)=(,d)`$ and to the sequence of maps $`g_k:=\gamma _kf_k`$, where $`\widehat{S}`$ is the one point compactification of $`S`$ and $`\widehat{\delta }`$ the distance on $`\widehat{S}`$ associated to $`\delta `$, cf. Lemma 2.2.
Let us denote $`x_0^k,x_1^k,x_2^k`$ be the points in $`S`$ such that $`f_k\left(x_i^k\right)=\theta _i^k`$, for $`i\{0,1,2\}`$.
Let us check that the assumptions of lemma 6.2 are verified.
The fact that $`lim_k\mathrm{}dist_H(\stackrel{~}{D}_k,\widehat{S})=0`$ comes the same way as in , (5.5).
By (6.3), we have, $`\delta (x_i^k,x_j^k)\frac{\lambda _k}{2}d(\theta _i^k,\theta _j^k)`$ and by assumption we then get
(7.5)
$$\delta (x_i^k,x_j^k)\frac{1}{2C}.$$
We then get the separation assumption on triples of points by choosing $`\delta :=inf\{D,\frac{1}{2C}\}`$.
It remains to check the assumption on $`g_k\left(\stackrel{~}{D}_k\right)=\gamma _kf_k\left(\stackrel{~}{D}_k\right)=\gamma _k\left(D_k\right)`$, namely,
(7.6)
$$lim_k\mathrm{}dist_H(\gamma _k\left(D_k\right),\mathrm{\Lambda }_C^{})=0.$$
In order to prove the property (6.6), we want to apply the lemma 6.3, but as the set $`(,d)`$ is a priori not uniformly perfect, we shall replace the uniform perfectness by the fact that $`(,\lambda _kd,\theta _0)`$ converges to a space $`(S,\delta _0,0)`$, which is uniformly perfect, cf. ().
We will show the
###### Lemma 7.4.
We consider the subsets $`\stackrel{~}{D}_kB_S(0,k)`$ and $`D_kB_{\lambda _k}(\theta _0,k)`$, where $`\stackrel{~}{D}_k`$ and $`D_k`$ are $`1/k`$-dense subsets of $`B_S(0,k)`$ and $`B_{(,\lambda _kd)}(\theta _0,k)`$, and the bijections $`f_k:\stackrel{~}{D}_kD_k`$ coming from the convergence of the sequence of pointed metric spaces $`(,\lambda _kd,\theta _0)`$ to $`(S,\delta ,0)`$ where $`(S,\delta ,0)`$ is a weak tangent of $`(\stackrel{~}{X},d)`$. We also assume that there exist positive constants $`C`$ and $`\delta `$, a sequence of points $`\theta _1^k,\theta _2^k\mathrm{\Lambda }_C^{}`$ and a sequence of elements $`\gamma _kC^{}`$ such that $`C^1\lambda _k^1d(\theta _i^k,\theta _j^k)C\lambda _k^1`$ and $`d(\gamma _k\theta _i^k,\gamma _k\theta _j^k)\delta `$ for all $`0ij2`$. Then, the Hausdorf distance $`dist_H(\gamma _kD_k,)`$ tends to $`0`$ as $`k`$ tends to infinity.
Proof : The proof is word by word the same as the proof of lemma 6.3, ie. lemma 5.1 (i) of with a difference in case 2).
We have $`B_{\lambda _k}(\theta _0,k)=B_{}(\theta _0,\frac{k}{\lambda _k})`$ and $`D_kB_{\lambda _k}(\theta _0,k)`$ an $`\frac{1}{k}`$-dense subset, for the metric $`\lambda _kd`$. In term of the distance $`d`$, the set $`D_k`$ is $`\left(ϵ_kR_k\right)`$-dense in $`B_{}(\theta _0,R_k)`$, where $`R_k:=\frac{k}{\lambda _k}`$ and $`ϵ_k:=\frac{1}{k^2}`$. By assumption, the points $`\theta _0^k=\theta _0,\theta _1^k,\theta _2^k`$ belong to $`B_{}(\theta _0,\mu _kR_k)`$, and satisfy
(7.7)
$$d(\theta _i^k,\theta _j^k)\delta _kR_k$$
where $`\delta _k:=\frac{1}{Ck}`$, and $`\mu _k:=\frac{C}{k}`$.
The points $`\gamma _k\theta _i^k`$ satisfy
(7.8)
$$d(\gamma _k\theta _i^k,\gamma _k\theta _j^k)\delta ,$$
and $`\frac{ϵ_k}{\delta _k^2}=C^2`$ is bounded.
Let us consider a point $`\theta `$. We want to approximate it by a point of $`\gamma _kD_k`$.
We can write $`\theta =\gamma _k\theta _k`$, for some $`\theta _k\mathrm{\Lambda }_C^{}`$. There are two cases.
Case 1). For infinitely many indices $`k`$, $`\theta _kB_{}(\theta _0,R_k)`$. We work in that case for these indices $`k`$, thus there are points $`\theta _k^{}D_kB_{}(\theta _0,R_k)`$, with $`d(\theta _k,\theta _k^{})ϵ_kR_k.`$
Since the distance between the $`\theta _i^k`$’s is bounded below by $`\delta _kR_k`$, we can find at least two of them which we call $`a_k`$ and $`b_k`$, such that
$$d(\theta _k,b_k)\frac{\delta _kR_k}{2}$$
and,
(7.9)
$$d(\theta _k^{},a_k)\frac{\delta _kR_k}{2}.$$
As $`C^{}`$ is contained in the cocompact group $`\mathrm{\Gamma }`$, it acts in a quasi-Möbius way on $`(\stackrel{~}{X},d)`$
thus,
(7.10)
$$\frac{d(\gamma _k\theta _k^{},\gamma _k\theta _k)d(\gamma _ka_k,\gamma _kb_k)}{d(\gamma _k\theta _k^{},\gamma _kb_k)d(\gamma _ka_k,\gamma _k\theta _k)}\eta \left(\frac{d(\theta _k^{},\theta _k)d(a_k,b_k)}{d(\theta _k^{},b_k)d(\theta _k,a_k)}\right)$$
for some homeomorphism $`\eta :[0,\mathrm{})[0,\mathrm{})`$. This implies
(7.11)
$$d(\gamma _k\theta _k^{},\gamma _k\theta _k)\frac{\left(diam\right)^2\eta \left(8ϵ_k\mu _k/\delta _k^2\right)}{\delta },$$
therefore $`d(\gamma _k\theta _k^{},\gamma _k\theta _k)`$ tends to zero as $`k`$ tends to infinity.
Case 2) . For all but finitely many indices $`k`$, $`\theta _kB_{}(\theta _0,R_k)`$.
We work with these indices $`k`$ such that $`\theta _kB_{}(\theta _0,R_k)`$.
We know that $`ϵ_k/\delta _k^2`$ is bounded above independantly of $`k`$, and by assumption, $`\delta _k2\mu _k`$.
We claim that there exist $`\xi _kD_k`$ and a positive constant $`c_0`$ such that for all $`k`$,
(7.12)
$$\frac{d(\xi _k,\theta _0)}{R_k}c_0$$
let us prove the claim.
On one hand, as $`(\stackrel{~}{X},d)`$ is uniformly perfect, and so is it’s weak tangent $`(S,\delta )`$ because the one point compactification $`(\widehat{S},\widehat{\delta })`$ of $`(S,\delta )`$ is quasi-Möbius homeomorphic to $`(\stackrel{~}{X},d)`$, therefore there exist a constant $`C_0[0,1)`$ such that for every $`xS`$ and $`0<R<diamS`$, we have
(7.13)
$$\overline{B}_{(S,\delta )}(0,R)B_{(S,\delta )}(0,C_0R)\mathrm{}.$$
On the other hand, $`(,\lambda _kd,\theta _0)`$ converges to $`(S,\delta ,0)`$. After reindexing the sequence $`\left\{\lambda _k\right\}`$, we have for each $`ϵ>0`$ a map $`g_k:B_{\lambda _k}(\theta _0,k)S`$ such that
(i) $`g_k\left(\theta _0\right)=0`$,
for any two points $`\theta `$ and $`\theta ^{}`$ in $`B_{\lambda _k}(\theta _0,k)`$,
(ii) $`\left|\delta (g_k\left(\theta \right),g_k\left(\theta ^{}\right))\lambda _kd(\theta ,\theta ^{})\right|ϵ`$,
(iii) the $`ϵ`$-neighborhood of $`g_k\left(B_{\lambda _k}(\theta _0,k)\right)`$ contains $`B_{(S,\delta )}(0,kϵ)`$.
By (iii), we have
(7.14)
$$\overline{B}_{(S,\delta )}(0,kϵ)𝒰_ϵ^{(S,\delta )}g_k\left(\overline{B}_{\lambda _k}(\theta _0,k)\right).$$
By (6.13) there exist $`y_k\overline{B}_{(S,\delta )}(0,kϵ)B_{(S,\delta )}(0,C_0\left(kϵ\right))`$, and by (6.14) there exist $`\xi _k^{}\overline{B}_{\lambda _k}(\theta _0,k)`$ such that
(7.15)
$$\delta (y_k,g_k\left(\xi _k^{}\right))ϵ.$$
We now evaluate $`d(y_k,g_k\left(\xi _k^{}\right))`$. By the above properties (i), (ii), (6.15) and the triangle inequality we have
$$\lambda _kd(\xi _k^{},\theta _0)\delta (g_k\left(\xi _k^{}\right),0)ϵ\delta (y_k,0)\delta (y_k,g_k\left(\xi _k^{}\right))ϵ$$
(7.16)
$$C_0\left(kϵ\right)2ϵ.$$
As $`D_k`$ is $`ϵ_kR_k`$-dense in $`B_{(,d)}(\theta _0,k/\lambda _k)`$, there exist $`\xi _kD_k`$ such that $`d(\xi _k,\xi _k^{})ϵ_kR_k=\frac{kϵ_k}{\lambda _k}`$.
Let us denote $`c_0=C_0/2`$. For $`k`$ large enough we have $`\frac{C_0\left(kϵ\right)2ϵkϵ_k}{\lambda _k}\frac{c_0k}{\lambda _k}`$, therefore by (6.16) we get
(7.17)
$$d(\xi _k,\theta _0)d(\xi _k^{},\theta _0)d(\xi _k^{},\xi _k)c_0R_k,$$
which proves the claim.
We can assume that for $`k`$ large enough, $`\mu _k<c_0/2<1/2`$.
We choose $`a_k=\theta _1^k`$ and $`b_k=\theta _2^k`$, and we get
$$\frac{d(\gamma _k\xi _k,\gamma _k\theta _k)d(\gamma _ka_k,\gamma _kb_k)}{d(\gamma _k\xi _k,\gamma _kb_k)d(\gamma _ka_k,\gamma _k\theta _k)}\eta \left(\frac{d(\xi _k,\theta _k)d(a_k,b_k)}{d(\xi _k,b_k)d(\theta _k,a_k)}\right)$$
$$\eta \left(\frac{4\mu _kd(\theta _k,\theta _0^k)}{(d(\theta _k,\theta _0^k)\mu _kR_k)\left)\right(c_0\mu _k)}\right)$$
(7.18)
$$\eta \left(16\mu _k/c_0\right).$$
We get
$$d(\gamma _k\theta _k,\gamma _k\xi _k)\left(diam\mathrm{\Lambda }_C^{}\right)^2\eta \left(16\mu _k/c_0\right)/\delta .$$
$`\mathrm{}`$
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# Brief Report A simple scheme of teleportation of arbitrary multipartite qubit entanglement
## Abstract
In this paper, we define a cross product operator and construct the cross Bell basis, by use this basis and Bell measurements we give a simple scheme of the teleportation of arbitrary multipartite qubit entanglement.
PACC numbers: 03.67.Mn, 03.65.Ud, 03.67.Hk.
In modern quantum mechanics and quantum information, the quantum teleportation is a quite interesting and important topic. Following the BBCJPW scheme, there have been very many related works (e.g. see the references in ). Generally, these works discussed the teleportaion of the single unknown qubit states. Recently, Rigolin gives a schemes of the teleportation of arbitrary multipartite qubit states, there are yet some relational improvements (e.g., see ). However these schemes are more complex whatever in physics and in mathematical forms. Generally, in a perfect quantum teleportation scheme, the most basic matter are to find the quantum channels (they best form a basis of the a Hilbert space) and to give a real physical measurement way distinguishing the outcomes of the wave function collapses. In this paper, we give a simple scheme of teleportation of arbitrary multipartite qubit entanglement. Our ways are to define a convenient operator ‘cross products’, to construct the ‘cross Bell basis’ and to use two or more common Bell measurements.
In the following we write the Hilbert space of states of a spin-$`\frac{1}{2}`$ particle $`x`$ as $`H_x`$, in which a pure-states $`\mathrm{\Psi }_x=\underset{i=0,1}{}c_ii_x.`$ In this paper, we mainly discuss the Hilbert space $`H_1H_2H_3H_4.`$
Definition. Suppose that $`\mathrm{\Psi }_{13}=\underset{i,j=0,1}{}c_{ij}i_1j_3H_1H_3,`$ $`\mathrm{\Phi }_{24}=\underset{r,s=0,1}{}d_{rs}r_2s_4H_2H_4`$ are two pure-states, then the cross product $`\mathrm{\Psi }_{13}\mathrm{}\mathrm{\Phi }_{24}H_1H_2H_3H_4`$ of $`\mathrm{\Psi }_{13}`$ and $`\mathrm{\Phi }_{24}`$ is defined to be the result of $`\mathrm{\Psi }_{13}\mathrm{\Phi }_{24}`$ returning to the natural order 1,2,3,4, i.e.
$$\mathrm{\Psi }_{13}\mathrm{}\mathrm{\Phi }_{24}=\underset{i,r,j,s=0,1}{}c_{ij}d_{rs}i_1r_2j_3s_4$$
(1)
Notice that since the order of $`H_1,H_2,H_3,H_4`$ is important in our discussion, so the operations $``$ and $`\mathrm{}`$ are distinct. In addition, obviously $`\mathrm{}`$ is a bi-linear and non-commutative operator.
Now, we read the ordinary Bell bases as
$$\mathrm{\Psi }_{\alpha \beta }^\pm =\frac{1}{\sqrt{2}}\left(0_\alpha 0_\beta \pm 1_\alpha 1_\beta \right),\mathrm{\Phi }_{\alpha \beta }^\pm =\frac{1}{\sqrt{2}}\left(0_\alpha 1_\beta \pm 1_\alpha 0_\beta \right)$$
(2)
then by using crossed products and according to the rule similar to matrix entries, we can write a group $`𝔹=\left\{K\mathrm{}L\right\}`$ of sixteen pure-states $`K\mathrm{}L`$ as
$$𝔹=\begin{array}{ccccc}^{_{K\mathrm{}L}}\hfill & ^{_{_{\mathrm{\Psi }_{24}^+}}}\hfill & ^{_{_{\mathrm{\Psi }_{24}^{}}}}\hfill & ^{_{^{_{\mathrm{\Phi }_{24}^+}}}}\hfill & ^{_{_{\mathrm{\Phi }_{24}^{}}}}\hfill \\ ^{_{\mathrm{\Psi }_{13}^+}}\hfill & \mathrm{\Psi }_{13}^+\mathrm{}\mathrm{\Psi }_{24}^+,\hfill & \mathrm{\Psi }_{13}^+\mathrm{}\mathrm{\Psi }_{24}^{},\hfill & \mathrm{\Psi }_{13}^+\mathrm{}\mathrm{\Phi }_{24}^+,\hfill & \mathrm{\Psi }_{13}^+\mathrm{}\mathrm{\Phi }_{24}^{}\hfill \\ ^{_{\mathrm{\Psi }_{13}^{}}}\hfill & \mathrm{\Psi }_{13}^{}\mathrm{}\mathrm{\Psi }_{24}^+,\hfill & \mathrm{\Psi }_{13}^{}\mathrm{}\mathrm{\Psi }_{24}^{},\hfill & \mathrm{\Psi }_{13}^{}\mathrm{}\mathrm{\Phi }_{24}^+,\hfill & \mathrm{\Psi }_{13}^{}\mathrm{}\mathrm{\Phi }_{24}^{}\hfill \\ ^{_{\mathrm{\Phi }_{13}^+}}\hfill & \mathrm{\Phi }_{13}^+\mathrm{}\mathrm{\Psi }_{24}^+,\hfill & \mathrm{\Phi }_{13}^+\mathrm{}\mathrm{\Psi }_{24}^{},\hfill & \mathrm{\Phi }_{13}^+\mathrm{}\mathrm{\Phi }_{24}^+,\hfill & \mathrm{\Phi }_{13}^+\mathrm{}\mathrm{\Phi }_{24}^{}\hfill \\ ^{_{\mathrm{\Phi }_{13}^{}}}\hfill & \mathrm{\Phi }_{13}^{}\mathrm{}\mathrm{\Psi }_{24}^+,\hfill & \mathrm{\Phi }_{13}^{}\mathrm{}\mathrm{\Psi }_{24}^{},\hfill & \mathrm{\Phi }_{13}^{}\mathrm{}\mathrm{\Phi }_{24}^+,\hfill & \mathrm{\Phi }_{13}^{}\mathrm{}\mathrm{\Phi }_{24}^{}\hfill \end{array}$$
(3)
It is easily verified that $`𝔹`$ is a complete orthogonal basis of $`H_1H_2H_3H_4`$, we call it the ‘crossed Bell basis’. Here it must be stressed that these bases are really distinguishable by Bell measurements. For instance, for any state $`\mathrm{\Psi }_{1234}H_1H_2H_3H_4`$ if we make two independent Bell measurements jointed particle pair $`(1,3)`$ and jointed particle pair $`(2,4)`$ respectively, then $`\mathrm{\Psi }_{1234}`$ must collapse to one of the above sixteen crossed Bell bases with a probability. In the following, we notice $`\stackrel{}{x}=1x`$ for $`x=0`$ or $`1.`$ The transformation from the natural basis to the crossed Bell basis is
$``$ $`i_1r_2i_3r_4={\displaystyle \frac{1}{2}}(\mathrm{\Psi }_{13}^++\mathrm{\Psi }_{13}^{})\mathrm{}(\mathrm{\Psi }_{24}^++\mathrm{\Psi }_{24}^{})`$ (4)
$``$ $`i_1r_2i_3\underset{4}{\overset{}{r}}={\displaystyle \frac{1}{2}}(1)^r(\mathrm{\Psi }_{13}^++\mathrm{\Psi }_{13}^{})\mathrm{}(\mathrm{\Phi }_{24}^++\mathrm{\Phi }_{24}^{})`$ (5)
$``$ $`i_1r_2\underset{3}{\overset{}{i}}r_4={\displaystyle \frac{1}{2}}(1)^i(\mathrm{\Phi }_{13}^++\mathrm{\Phi }_{13}^{})\mathrm{}(\mathrm{\Psi }_{24}^++\mathrm{\Psi }_{24}^{})`$ (6)
$``$ $`i_1r_2\underset{3}{\overset{}{i}}\underset{4}{\overset{}{r}}={\displaystyle \frac{1}{2}}(1)^{i+r}(\mathrm{\Phi }_{13}^++\mathrm{\Phi }_{13}^{})\mathrm{}(\mathrm{\Phi }_{24}^++\mathrm{\Phi }_{24}^{})`$ (7)
Now we prove that by using any one of cross Bell bases, we can realize the teleportation of a unknown bipartite qubit pure-state. For instance, we take $`\mathrm{\Phi }_{13}^+`$ $`\mathrm{}\mathrm{\Phi }_{24}^{}`$, as the quantum channel, particles 3, 4 are in Alice. The receiptor is Bob, she is in remote place from Alice, and she holds particles 1, 2. Suppose that $`\phi _{56}=\alpha 0_50_6+\beta 0_51_6+\gamma 1_50_6+\delta 1_51_6`$ is a client unknown state in Alice. It is known that $`\psi _{56}`$ is entangled if and only if $`\alpha \gamma \beta \delta `$ $`0`$.
In the present case the total state is
$``$ $`\mathrm{\Psi }_{123456}=(\mathrm{\Phi }_{13}^+\mathrm{}\mathrm{\Phi }_{24}^{})\phi _{56}`$ (8)
$`=`$ $`{\displaystyle \frac{1}{2}}\left(0_10_21_31_40_11_21_30_4+1_10_20_31_41_11_20_30_4\right)`$ (10)
$`\left(\alpha 0_50_6+\beta 0_51_6+\gamma 1_50_6+\delta 1_51_6\right)`$
Expanding $`\mathrm{\Psi }_{123456}`$ and by using the formulae in Eq.(2) to particles 3, 4, 5 and 6, the final result is
$$\mathrm{\Psi }_{123456}=\left\{\begin{array}{c}\frac{1}{4}\left(\delta 0_10_2+\gamma 0_11_2\beta 1_10_2\alpha 1_11_2\right)\mathrm{\Psi }_{35}^+\mathrm{}\mathrm{\Psi }_{46}^+\\ +\frac{1}{4}\left(\delta 0_10_2+\gamma 0_11_2+\beta 1_10_2\alpha 1_11_2\right)\mathrm{\Psi }_{35}^+\mathrm{}\mathrm{\Psi }_{46}^{}\\ +\frac{1}{4}\left(\gamma 0_10_2+\delta 0_11_2\alpha 1_10_2\beta 1_11_2\right)\mathrm{\Psi }_{35}^+\mathrm{}\mathrm{\Phi }_{46}^+\\ +\frac{1}{4}\left(\gamma 0_10_2+\delta 0_11_2+\alpha 1_10_2\beta 1_11_2\right)\mathrm{\Psi }_{35}^+\mathrm{}\mathrm{\Phi }_{46}^{}\\ \\ +\frac{1}{4}\left(\delta 0_10_2\gamma 0_11_2\beta 1_10_2\alpha 1_11_2\right)\mathrm{\Psi }_{35}^{}\mathrm{}\mathrm{\Psi }_{46}^+\\ +\frac{1}{4}\left(\delta 0_10_2\gamma 0_11_2+\beta 1_10_2\alpha 1_11_2\right)\mathrm{\Psi }_{35}^{}\mathrm{}\mathrm{\Psi }_{46}^{}\\ +\frac{1}{4}\left(\gamma 0_10_2\delta 0_11_2\alpha 1_10_2\beta 1_11_2\right)\mathrm{\Psi }_{35}^{}\mathrm{}\mathrm{\Phi }_{46}^+\\ +\frac{1}{4}\left(\gamma 0_10_2\delta 0_11_2+\alpha 1_10_2\beta 1_11_2\right)\mathrm{\Psi }_{35}^{}\mathrm{}\mathrm{\Phi }_{46}^+\\ \\ +\frac{1}{4}\left(\beta 0_10_2+\alpha 0_11_2\delta 1_10_2\gamma 1_11_2\right)\mathrm{\Phi }_{35}^+\mathrm{}\mathrm{\Psi }_{46}^+\\ +\frac{1}{4}\left(\beta 0_10_2+\alpha 0_11_2+\delta 1_10_2\gamma 1_11_2\right)\mathrm{\Phi }_{35}^+\mathrm{}\mathrm{\Psi }_{46}^{}\\ +\frac{1}{4}\left(\alpha 0_10_2+\beta 0_11_2\gamma 1_10_2\delta 1_11_2\right)\mathrm{\Phi }_{35}^+\mathrm{}\mathrm{\Phi }_{46}^+\\ +\frac{1}{4}\left(\alpha 0_10_2+\beta 0_11_2+\gamma 1_10_2\delta 1_11_2\right)\mathrm{\Phi }_{35}^+\mathrm{}\mathrm{\Phi }_{46}^{}\\ \\ +\frac{1}{4}\left(\beta 0_10_2\alpha 0_11_2\delta 1_10_2\gamma 1_11_2\right)\mathrm{\Phi }_{35}^{}\mathrm{}\mathrm{\Psi }_{46}^+\\ +\frac{1}{4}\left(\beta 0_10_2\alpha 0_11_2+\delta 1_10_2\gamma 1_11_2\right)\mathrm{\Phi }_{35}^{}\mathrm{}\mathrm{\Psi }_{46}^{}\\ +\frac{1}{4}\left(\alpha 0_10_2\beta 0_11_2\gamma 1_10_2\delta 1_11_2\right)\mathrm{\Phi }_{35}^{}\mathrm{}\mathrm{\Phi }_{46}^+\\ \left(\alpha 0_10_2\beta 0_11_2+\gamma 1_10_2\delta 1_11_2\right)\mathrm{\Phi }_{35}^{}\mathrm{}\mathrm{\Phi }_{46}^{}\end{array}\right\}$$
(11)
We define eight $`2\times 2`$ unitary matrices by
$`U_{\mathrm{\Psi }_{35}^+}`$ $`=`$ $`i\sigma _y,U_{\mathrm{\Psi }_{35}^{}}=\sigma _x,U_{\mathrm{\Phi }_{35}^+}=\sigma _z,U_{\mathrm{\Phi }_{35}^{}}=\sigma _0`$ (12)
$`U_{\mathrm{\Psi }_{46}^+}`$ $`=`$ $`\sigma _x,U_{\mathrm{\Psi }_{46}^{}}=i\sigma _y,U_{\mathrm{\Phi }_{46}^+}=\sigma _0,U_{\mathrm{\Phi }_{46}^{}}=\sigma _z`$ (13)
where $`\sigma _x=\left[\begin{array}{cc}& 1\hfill \\ 1\hfill & \end{array}\right],\sigma _y=\left[\begin{array}{cc}& i\hfill \\ i\hfill & \end{array}\right],\sigma _z=\left[\begin{array}{cc}1\hfill & \\ & 1\hfill \end{array}\right],\sigma _0=\left[\begin{array}{cc}1\hfill & \\ & 1\hfill \end{array}\right]`$ are the Pauli matrices. Now $`\mathrm{\Psi }_{123456}`$ can be written simply as
$``$ $`\mathrm{\Psi }_{123456}={\displaystyle _{K=\mathrm{\Psi }_{35}^+,\mathrm{\Psi }_{35}^{},\mathrm{\Phi }_{35}^+,\mathrm{\Phi }_{35}^+}}{\displaystyle _{L=\mathrm{\Psi }_{46}^+,\mathrm{\Psi }_{46}^{},\mathrm{\Phi }_{46}^+,\mathrm{\Phi }_{46}^+}}({\displaystyle \frac{1}{4}}\phi _{K\mathrm{}L}^{}\left(K\mathrm{}L\right))`$ (14)
$``$ $`\phi _{K\mathrm{}L}^{}=U_KU_L(\phi _{12})`$ (15)
where $`\phi _{12}=\phi _{56}\left(51,62\right)`$. When Alice makes a Bell measurement of particle pair $`(3,5)`$, and a Bell measurement of particle pair $`(4,6)`$ respectively, the wave function will collapse to one $`\phi _{K\mathrm{}L}^{}\left(K\mathrm{}L\right)`$ with probability $`\frac{1}{16}`$ ($`K\mathrm{}L`$ is measured by Alice, simultaneously Bob obtain a corresponding state $`\phi _{K\mathrm{}L}^{}).`$ When Alice informs Bob of her measurement result (one $`K\mathrm{}L)`$ by a classical communication, then Bob at once knows that the correct result should be
$$\phi _{12}=\left(U_KU_L\right)^1\left(\phi _{K\mathrm{}L}^{}\right)=U_L^TU_K^T\left(\phi _{K\mathrm{}L}^{}\right)$$
(16)
where $`T`$is the transposition. So, the bipartite qubit entanglement teleportation has been completed. If we take other cross Bell basis as channel, the steps are similar.
We see that in our method the order (13 and 24) of particles in cross Bell basis is important, in fact, if we use the natural order (12 and 34), e.g. $`\mathrm{\Psi }_{12}^\pm ,\mathrm{\Phi }_{34}^\pm ,\mathrm{},`$ etc., the product $`\mathrm{}`$ becomes common tensor product $``$, and we still choose a product of them as the quantum channel, then the process in the above scheme will lead to inconveniency and difficulty.
Discussion and conclusion. If it is known that when $`\alpha \gamma \beta \delta =0,`$ then $`\phi _{56}`$ must be decomposed in form as $`\phi _{56}=\left(a0_5+b1_5\right)\left(c0_6+d1_6\right)`$, then obviously the above process, in fact, becomes two independent teleportation of $`\phi _5=a0_5+b1_5`$ and $`\phi _6=c0_6+d1_6`$ respectively. In addition, for three states $`\mathrm{\Psi }_{14}=\underset{i,j=0,1}{}c_{ij}i_1j_4H_1H_4,`$ $`\mathrm{\Phi }_{25}=\underset{r,s=0,1}{}d_{rs}r_2s_5H_2H_5`$ and $`\mathrm{\Omega }_{36}=\underset{x,y=0,1}{}e_{xy}x_3y_6H_3H_6`$ if we define the cross product
$$\mathrm{\Psi }_{14}\mathrm{}\mathrm{\Phi }_{25}\mathrm{}\mathrm{\Omega }_{36}=\underset{i,r,j,s=0,1}{}c_{ij}d_{rs}e_{xy}i_1r_2x_3j_4s_5y_6\underset{m=1}{\overset{6}{}}H_m$$
(17)
and construct the cross Bell basis $`\left\{K\mathrm{}L\mathrm{}M\right\}`$ of $`\underset{m=1}{\overset{6}{}}H_m`$, where $`K,L,M=\mathrm{\Psi }_{14}^\pm ,\mathrm{\Phi }_{14}^\pm ,`$ $`L=\mathrm{\Psi }_{25}^\pm ,\mathrm{\Phi }_{25}^\pm `$,$`M=\mathrm{\Psi }_{36}^\pm ,\mathrm{\Phi }_{36}^\pm `$, etc., then by a similar way we can realize the teleporation of a unknown tripartite qubit state. Obviously this method can be generalized to arbitrary dimensional cases.
To sum up, by using cross Bell bases and Bell measurements we give a simple scheme of arbitrary multipartite qubit states.
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# Diffusive spreading and mixing of fluid monolayers
## 1 Introduction
There are substantial efforts to miniaturize chemical processes by using microfluidic systems. The “lab on a chip concept” integrates a great variety of chemical and physical processes into a single device in a similar way as an integrated circuit incorporates many electronic devices into a single chip . These microfluidic devices do not only allow for cheap mass production but they can operate with much smaller quantities of reactants and reaction products than standard laboratory equipments. This is particularly important for solutions containing rare and expensive substances, such as certain biological materials, and for toxic or explosive components . Even though most microfluidic devices available today have micron sized channels, further miniaturization is leading towards the nano-scale . Besides meeting technical challenges, new theoretical concepts are needed to understand the basic physical processes underlying this new technology . Whereas the ultimate limits for miniaturization of electronic devices are set by quantum fluctuations, in a chemical chip these limits are determined by thermal fluctuations and can be explored by methods of classical statistical mechanics.
At the sub-micron thickness scale, recent experiments of liquid spreading on atomically smooth surfaces , performed with volumes of the order of nano-liters, have clearly shown by means of dynamic ellipsometry or X-ray reflectivity measurements that precursor films with molecular thickness and macroscopic extent advance in front of the macroscopic liquid wedge of the spreading drop. ( Thin, i.e., of the order of 100 nm, precursor films spreading ahead of the macroscopic droplet have also been observed experimentally .) The occurrence of molecularly thin precursor films with a similar spreading dynamics has been evidenced very recently also for immiscible metal systems in three regimes: solid drops with a solid film, solid drops with a liquid film, and liquid drops with a liquid film . Theoretical work (see Refs. and references therein) combined with an impressive number of Molecular Dynamics (MD) and Monte Carlo (MC) simulations (see Ref. and references therein) addressed the mechanisms behind the extraction and the experimentally observed $`t^{1/2}`$ asymptotic time dependence of the linear extent of the precursor films on chemically homogeneous substrates . This led to a good understanding of the spreading dynamics and of the intrinsic morphology of the films. Based on these results, more complicated issues can be addressed such as, e.g., the spreading behavior of monolayers exposed to chemically patterned substrates , or the question of mixing of different fluids at the nano-scale that we shall present below.
The organization of the paper is as follows. In Sec. 2 we briefly present the lattice gas model of interacting particles and discuss the rules defining the microscopic dynamics and the nonlinear diffusion equation derived from it in the continuum limit. Section 3 is devoted to a qualitative discussion of the results obtained for the case of monolayer spreading in the presence of mesoscopic obstacles, with a particular focus on the relaxation of the density profile upon encountering and passing the obstacle. In Sec. 4 we discuss the mixing of two species during spreading of monolayers following the merger of two chemical lanes, and we conclude with a brief summary of the results in Sec. 5.
## 2 Fluid monolayers on homogeneous substrates
Recently, we have studied the structure of a monolayer, which extracts from a reservoir and spreads on a flat, chemically homogeneous substrate, by using a lattice gas model of interacting particles as proposed in Refs. . Since we shall use this model as a starting point for our present study, for clarity and further reference we briefly describe the defining rules of the model and the non-linear diffusion equation obtained from the microscopic dynamics within the continuum limit. A thorough analysis of this model is presented in Ref. .
(a) We choose a homogeneous substrate such that the spreading occurs in the $`xy`$ plane. The half-plane $`x<0`$ is occupied by a reservoir of particles at fixed chemical potential which maintains at its contact line with the substrate — positioned at the line $`x=0`$ — an average density $`C_0`$ (defined as the number of particles per unit length in the transversal $`y`$ direction). At time $`t=0`$, the half-plane $`x>0`$ is empty. There is no imposed flow of particles from the reservoir pushing the extracting film.
(b) The substrate-fluid interaction is modeled as a periodic potential forming a lattice of potential wells with coordination number $`z`$ ($`z=4`$ for a square lattice) and lattice constant $`a`$. The particle motion proceeds via activated jumps between nearest-neighbor wells; evaporation from the substrate is not allowed. The activation barrier $`U_A`$ determines the jumping rate $`\mathrm{\Omega }=\nu _0\mathrm{exp}[U_A/k_BT]`$, where $`\nu _0`$ is an attempt frequency defining the time unit, $`k_B`$ is the Boltzmann constant, and $`T`$ is the temperature.
(c) The pair interaction between fluid particles at distance $`r`$ is taken to be hard-core repulsive at short range, preventing double occupancy of the wells, and attractive at long range, $`U_0/r^6`$ for $`r1`$, resembling a Lennard-Jones type interaction potential. Here and in the following all distances are measured in units of the lattice constant $`a`$ and therefore are dimensionless. The selection of the nearest-neighbor well into which a particle attempts to jump, i.e., the probability $`p(\mathrm{r}\mathbf{}\mathrm{r}^{\mathbf{}}\mathbf{;}\mathrm{t}\mathbf{)}`$ that a jump from location $`\mathrm{r}`$ will be directed toward the location $`\mathrm{r}^{\mathbf{}}`$, is biased by the fluid-fluid energy landscape and is given by
$$p(\mathrm{r}\mathbf{}\mathrm{r}^{\mathbf{}}\mathbf{;}\mathrm{t}\mathbf{)}\mathbf{=}\frac{\mathrm{𝐞𝐱𝐩}\mathbf{\left\{}\frac{𝜷}{2}\mathbf{[}\stackrel{\mathbf{~}}{\mathrm{U}}\mathbf{(}\mathrm{r}\mathbf{;}\mathrm{t}\mathbf{)}\mathbf{}\stackrel{\mathbf{~}}{\mathrm{U}}\mathbf{(}\mathrm{r}^{\mathbf{}}\mathbf{;}\mathrm{t}\mathbf{)}\mathbf{]}\mathbf{\right\}}}{\mathrm{Z}\mathbf{(}\mathrm{r}\mathbf{;}\mathrm{t}\mathbf{)}}\mathbf{,}$$
(1)
where $`Z(\mathrm{r}\mathbf{;}\mathrm{t}\mathbf{)}\mathbf{=}{\displaystyle \underset{\mathrm{r}^{\mathbf{}}\mathbf{,}\mathbf{|}\mathrm{r}^{\mathbf{}}\mathbf{}\mathrm{r}\mathbf{|}\mathbf{=}1}{\mathbf{}}}\mathrm{𝐞𝐱𝐩}\mathbf{\left\{}{\displaystyle \frac{𝜷}{2}}\mathbf{[}\stackrel{\mathbf{~}}{\mathrm{U}}\mathbf{(}\mathrm{r}\mathbf{;}\mathrm{t}\mathbf{)}\mathbf{}\stackrel{\mathbf{~}}{\mathrm{U}}\mathbf{(}\mathrm{r}^{\mathbf{}}\mathbf{;}\mathrm{t}\mathbf{)}\mathbf{]}\mathbf{\right\}}`$ is the normalization constant and $`1/\beta =k_BT`$,
$$\stackrel{~}{U}(\mathrm{r}\mathbf{;}\mathrm{t}\mathbf{)}\mathbf{=}\mathbf{}\mathrm{U}_0\underset{\mathrm{r}^{\mathbf{}}\mathbf{,}0\mathbf{<}\mathbf{|}\mathrm{r}^{\mathbf{}}\mathbf{}\mathrm{r}\mathbf{|}\mathbf{}3}{\mathbf{}}\frac{𝜼\mathbf{(}\mathrm{r}^{\mathbf{}}\mathbf{;}\mathrm{t}\mathbf{)}}{\mathbf{|}\mathrm{r}\mathbf{}\mathrm{r}^{\mathbf{}}\mathbf{|}^6}\mathbf{,}$$
(2)
and $`\eta (\mathrm{r}^{\mathbf{}}\mathbf{;}\mathrm{t}\mathbf{)}\mathbf{}\mathbf{\{}0\mathbf{,}1\mathbf{\}}`$ is the occupation number of the well at $`\mathrm{r}^{\mathbf{}}`$ at the time $`t`$. The summation in Eq. (2) has been restricted to three lattice units for computational convenience. This corresponds to the cut-off generally used in Molecular Dynamics simulations for algebraically decaying Lennard-Jones pair-potentials. The rates
$$\omega _{\mathrm{r}\mathbf{}\mathrm{r}^{\mathbf{}}\mathbf{;}\mathrm{t}}=\mathrm{\Omega }p(\mathrm{r}\mathbf{}\mathrm{r}^{\mathbf{}}\mathbf{;}\mathrm{t}\mathbf{)}$$
(3)
for the transitions from $`\mathrm{r}`$ to neighboring sites $`\mathrm{r}^{\mathbf{}}`$ satisfy
$$\underset{\mathrm{r}^{\mathbf{}}\mathbf{,}\mathbf{|}\mathrm{r}^{\mathbf{}}\mathbf{}\mathrm{r}\mathbf{|}\mathbf{=}1}{}\omega _{\mathrm{r}\mathbf{}\mathrm{r}^{\mathbf{}}\mathbf{;}\mathrm{t}}\mathrm{\Omega }.$$
(4)
Thus for any given particle at any location the total rate of leaving a potential well is determined only by the fluid-solid interaction characterized by $`U_A`$, is time-independent, and equals $`\mathrm{\Omega }`$.
Neglecting all spatial and temporal correlations, i.e., assuming that averages of products of occupation numbers $`\eta (\mathrm{r}\mathbf{;}\mathrm{t}\mathbf{)}`$ are equal to the corresponding products of averaged occupation numbers $`\rho (\mathrm{r}\mathbf{;}\mathrm{t}\mathbf{)}\mathbf{=}\mathbf{}𝜼\mathbf{(}\mathrm{r}\mathbf{;}\mathrm{t}\mathbf{)}\mathbf{}`$, where $`\mathrm{}`$ denotes the average with respect to the corresponding probability distribution $`𝒫(\{\eta (\mathrm{r}\mathbf{;}\mathrm{t}\mathbf{)}\mathbf{\}}\mathbf{)}`$ of a configuration $`\{\eta (\mathrm{r}\mathbf{;}\mathrm{t}\mathbf{)}\mathbf{\}}`$, one can formulate a mean-field master equation for the local occupational probability, i.e., the number density $`\rho (\mathrm{r}\mathbf{;}\mathrm{t}\mathbf{)}`$ . In the continuum limit of space and time ($`\mathrm{\Delta }t0`$, $`a0`$, $`\mathrm{\Omega }^10`$, $`D_0=\mathrm{\Omega }a^2/4`$ finite) for the master equation, by taking Taylor expansions for $`p(\mathrm{r}\mathbf{}\mathrm{r}^{\mathbf{}}\mathbf{)}`$ and $`\rho (\mathrm{r}^{\mathbf{}}\mathbf{;}\mathrm{t}\mathbf{)}`$ around $`\mathrm{r}`$ and keeping terms up to second-order spatial derivatives of the density $`\rho (\mathrm{r}\mathbf{;}\mathrm{t}\mathbf{)}`$ , one obtains the following nonlinear and nonlocal equation for $`\rho (\mathrm{r}\mathbf{;}\mathrm{t}\mathbf{)}`$ :
$$_t\rho =D_0\left[\rho +\beta \rho (1\rho )U\right]+𝒪(a^2)$$
(5)
where
$$U(\mathrm{r}\mathbf{;}\mathrm{t}\mathbf{)}\mathbf{}\mathbf{}\stackrel{\mathbf{~}}{\mathrm{U}}\mathbf{(}\mathrm{r}\mathbf{;}\mathrm{t}\mathbf{)}\mathbf{}\mathbf{=}\mathbf{}\mathrm{U}_0\underset{\mathrm{r}^{\mathbf{\prime \prime }}\mathbf{,}\mathrm{\hspace{0.17em}0}\mathbf{<}\mathbf{|}\mathrm{r}^{\mathbf{\prime \prime }}\mathbf{}\mathrm{r}\mathbf{|}\mathbf{}3}{\mathbf{}}\frac{𝝆\mathbf{(}\mathrm{r}^{\mathbf{\prime \prime }}\mathbf{;}\mathrm{t}\mathbf{)}}{\mathbf{|}\mathrm{r}^{\mathbf{\prime \prime }}\mathbf{}\mathrm{r}\mathbf{|}^6}$$
(6)
is replacing $`\stackrel{~}{U}(\mathrm{r}\mathbf{;}\mathrm{t}\mathbf{)}`$ in the definition (1) for $`p(\mathrm{r}\mathbf{}\mathrm{r}^{\mathbf{}}\mathbf{)}`$.
Being nonlinear and, due to the term involving the interaction potential $`U(\mathrm{r}\mathbf{;}\mathrm{t}\mathbf{)}`$, nonlocal, Eq. (5) cannot be solved analytically and in most of the cases even the computation of a numerical solution is a difficult task. However, assuming that the density $`\rho (\mathrm{r}\mathbf{;}\mathrm{t}\mathbf{)}`$ is a slowly varying function of the spatial coordinates the potential $`U(\mathrm{r}\mathbf{;}\mathrm{t}\mathbf{)}`$ may be expanded as
$`U(\mathrm{r}\mathbf{;}\mathrm{t}\mathbf{)}`$ $`=U_0{\displaystyle \underset{\mathrm{r}^{\mathbf{}}\mathbf{,}\mathrm{\hspace{0.17em}0}\mathbf{<}\mathbf{|}\mathrm{r}^{\mathbf{}}\mathbf{}\mathrm{r}\mathbf{|}\mathbf{}3}{}}{\displaystyle \frac{\rho (\mathrm{r}^{\mathbf{}}\mathbf{;}\mathrm{t}\mathbf{)}}{|\mathrm{r}^{\mathbf{}}\mathbf{}\mathrm{r}\mathbf{|}^6}}`$ (7)
$`U_0\rho (\mathrm{r}\mathbf{;}\mathrm{t}\mathbf{)}{\displaystyle \underset{\mathrm{r}^{\mathbf{}}\mathbf{,}\mathrm{\hspace{0.17em}0}\mathbf{<}\mathbf{|}\mathrm{r}^{\mathbf{}}\mathbf{}\mathrm{r}\mathbf{|}\mathbf{}3}{\mathbf{}}}{\displaystyle \frac{1}{\mathbf{|}\mathrm{r}^{\mathbf{}}\mathbf{}\mathrm{r}\mathbf{|}^6}}\mathbf{+}𝓞\mathbf{(}\mathrm{a}^2\mathbf{)}\mathbf{,}`$
which leads to the local equation
$$_t\rho =D_0\{\left[1gW_0\rho (1\rho )\right]\rho \}+𝒪(a^2),$$
(8)
where $`W_0=\beta U_0`$, and $`g={\displaystyle \underset{1|\mathrm{r}\mathbf{|}\mathbf{}\mathrm{r}_\mathrm{c}}{}}|\mathrm{r}\mathbf{|}^\mathbf{}6`$ is a geometrical factor depending on the lattice type (e.g., square, triangular, etc.) and on the cut-off range of the potential. For the present case of a square lattice and a cut-off at $`r_c=3`$ one has $`g4.64`$.
Rescaling time as $`t\tau =D_0t`$ and defining an effective diffusion coefficient
$$D_e(\rho )=1gW_0\rho (1\rho ),$$
(9)
Eq. (8) may be written in the usual form of a diffusion equation:
$$_\tau \rho =\left[D_e(\rho )\rho \right]+𝒪(a^2).$$
(10)
The functional form of $`D_e(\rho )`$ (Eq. (9)) implies that for $`W_0>4/g`$ there will be values $`\rho _i`$ of the density for which $`D_e(\rho _i)<0`$. For parameters such that $`W_0<4/g`$, Eq. (10) is a proper diffusion equation (though non-linear), while for $`W_0>4/g`$ instabilities are expected in the range of densities for which $`D_e(\rho _i)<0`$, i.e., for $`\rho _i(\rho _\alpha ^{},\rho _\alpha ^+)`$ where
$$\rho _\alpha ^\pm =\frac{1}{2}\left(1\pm \sqrt{1\frac{4}{gW_0}}\right).$$
(11)
It is known that these instabilities lead to discontinuities in the density profile (“shocks”), i.e., they correspond to the formation of sharp interfaces. For the model defined by the rules (a)-(d), the value for the threshold interaction strength for which such interfaces emerge is predicted by the continuum theory as $`W_0^{(t)}=4/g0.86`$, which is significantly smaller than the lower bound estimate $`W_0^{(t)}>1`$ from KMC simulations. We attribute this to the mean-field character of the derivation of the continuum equation. Therefore it is necessary to include particle-particle correlations into the mean-field description. Since the dynamics is possible only by jumps into empty sites, one can argue that for $`z=4`$ the summation in $`g`$ should include at most three contributions from nearest neighbor sites. This leads to $`g3.64`$ and an estimate for the threshold interaction $`W_0^{(t)}1.1`$, in good agreement with the KMC results. For the rest of the analysis we shall use this corrected value of $`g`$. Additional support for this corrected value is provided by the analysis of the density profiles .
The constraint, that the reservoir keeps the mean density at $`x=0`$ at a fixed value $`C_0`$, implies the boundary condition
$$\rho (x=0,y;t)=C_0.$$
(12)
Depending on the particular system under study, additional boundary conditions may have to be satisfied. For example, for the case studied in Ref. , in the absence of formation of interfaces, i.e., for interactions $`W_0<W_0^{(t)}`$ and for large times, the density on the advancing edge $`X(t)`$ could be considered as fixed and equal to $`C_1`$,
$$\rho (x=X(t),y;t)=C_1,$$
(13)
where $`C_1=0.11`$ as inferred from the kinetic Monte Carlo (KMC) simulations. (This boundary condition (Eq. (13)) naturally occurred also in the theory of Burlatsky et al .) For that system, the absence of boundaries along the $`y`$direction and the $`y`$-independence of the boundary conditions at $`x=0`$ and $`x=X(t)`$ leads to an effectively one-dimensional problem and to a scaling solution $`\rho (x,t)=\stackrel{~}{C}(\lambda =x/\sqrt{t})`$. The analysis of Eq. (10) depends on whether $`W_0<W_0^{(t)}`$ or $`W_0>W_0^{(t)}`$. As shown in Ref. , in both cases the solutions are in excellent agreement with those obtained from KMC simulations; typical results are shown in Fig. 1.
While for the liquid-on-solid systems mentioned in the Introduction these intrinsic density profiles have not been measured yet, data of such density profiles are available for the immiscible metal systems studied in Ref. and they are in at least good qualitative agreement with the theoretical ones. Since the present model appears to provide a simple but realistic description of a fluid monolayer spreading on a homogeneous substrate, it is natural to use it as a starting point to address more complex problems, such as the spreading of monolayers on designed chemically heterogeneous substrates, or the mixing of monolayers.
## 3 Diffusive spreading around mesoscopic obstacles
In order to apply the model described in Sec. 2 to study the spreading of a monolayer around a mesoscopic-size obstacle, we add to the rules (a)-(d) in Sec. II the following ones:
(e) The obstacle is taken to be a square-shaped domain $`𝒟`$ of side length $`h`$ centered at $`(x=dh/2,y=0)`$ \[see also Fig. 2(a)\]. This domain is composed of sites with very low affinity for the fluid particles. The activation barrier $`U_𝒟`$ for jumps from sites outside $`𝒟`$ to those inside $`𝒟`$ is taken to be much larger than $`U_A`$, such that the boundary $`𝒟`$ of $`𝒟`$ acts effectively as a hard wall.
(f) A sink reservoir occupies the region $`xL`$, where $`L1`$ and $`Ld+h`$, and maintains at its contact line with the substrate, positioned at the line $`x=L`$, an average density (number of particles per unit length in the transversal $`y`$ direction) $`C_1=0`$.
Under these assumptions, the density profile $`\rho (x,y,t)`$ as the solution of Eq. (10) fulfills the initial condition
$$\rho (x,y,0)=C_0\mathrm{\Theta }(x),$$
(14)
where $`\mathrm{\Theta }(x)`$ denotes the Heaviside step function, and the boundary conditions
$`\rho (0,y,t)`$ $`=`$ $`C_0,`$
$`\rho (L,y,t)`$ $`=`$ $`0,`$ (15)
$`\mathrm{j}_\mathrm{n}\mathbf{|}{}_{\mathbf{}𝓓}{}^{}\mathbf{=}0\mathbf{.}`$
The current $`\mathrm{j}`$ is given by \[see Eq. (10)\]
$$\mathrm{j}\mathbf{=}\mathbf{}\mathrm{D}_\mathrm{e}\mathbf{(}𝝆\mathbf{)}\mathbf{}𝝆\mathbf{.}$$
(16)
Note that in a numerical study the system necessarily has a finite size $`L_y`$ along the $`y`$ direction. We will use periodic boundary conditions and sizes $`L_y1`$ such that $`h/L_y<1`$ but not negligible (mesoscopic-size obstacle) and $`L_yhr_c`$, such that the boundary $`𝒟`$ of the obstacle is sufficiently far away from the edge of the simulation box to avoid finite-size effects.
In Figs. 2(b) and (c) we present typical results for the density profiles in the vicinity of the obstacle obtained from the numerical integration of Eq. (10) with initial and boundary conditions given by Eqs. (14)-(3), respectively, for a spreading monolayer whose edge just encounters an obstacle ($`\tau =200`$)(b) and has just passed the obstacle ($`\tau =2000`$)(c). Several conclusions can be drawn from visually inspecting Figs. 2(b) and (c). Upon approaching the obstacle, the boundary condition of zero normal current at the boundary of the obstacle (reflecting wall) leads to an increase in the density in a region at the front of the spreading monolayer, as shown by the forward bending of the iso-density lines \[see, e.g., the yellow band in Fig. 2(b)\]. Upon passing the front-edge of the obstacle, the iso-density lines become straight, and once they reach the end of the obstructed region they bend again, this time backwards. This indicates that upon passing the obstacle the spreading tends to proceed faster in the regions far from the obstacle, while at the obstacle the iso-density lines are pinned until they cover the whole edge on the back of the obstacle \[see, e.g., the light blue band neighboring the green region in Fig. 2(c)\]. Once this is realized, the iso-density line detaches from the back-edge of the obstacle, and the bending slowly relaxes \[see, e.g., the boundary between the light and dark blue regions in the top region of Fig. 2(c)\], while the spreading continues; far away from the obstacle, the iso-density lines become again straight.
We end this section by noting that the maximum linear extent of the region where the iso-density lines are deformed, as well as the survival time of these deformations can be used as quantitative measures to describe the relaxation of the density perturbations induced by the obstacle as a function of the inter-particle attractive interactions, as well as of the scaled size $`h/L_y`$ of the mesoscopic obstacle (assuming that this is the most relevant geometrical parameter). The results of this analysis will be the subject of a forthcoming paper.
## 4 Diffusive mixing of two fluid monolayers composed of different species
The model described in Sec. 2 can be used to study the mixing of two spreading fluid monolayers composed of different species $`A`$ and $`B`$, if one assumes that the two species interact with the substrate in such a way that the same lattice structure of potential wells, eventually with different depths (i.e., different escape rates $`\mathrm{\Omega }_i`$), can accommodate both types of species. Assuming a square lattice of lattice constant $`a`$ and assuming the on-site hard core repulsion between any two particles such that double occupancy remains forbidden, the equations satisfied by the densities $`\rho _j(x,t)`$, where $`j\{A,B\}`$, are obtained from Eq. (5) by replacing $`U`$ with $`U_{ii}+U_{ij}`$, where $`U_{ii}`$ is the potential due to same-species interactions while $`U_{ij}`$, $`ji`$, is the potential due to interactions between different species, $`D_0`$ with $`D_i=\mathrm{\Omega }_ia^2/4`$, and changing the single-occupancy term from $`1\rho (\mathrm{r}\mathbf{,}\mathrm{t}\mathbf{)}`$ to $`1\rho _A(\mathrm{r}\mathbf{,}\mathrm{t}\mathbf{)}𝝆_\mathrm{B}\mathbf{(}\mathrm{r}\mathbf{,}\mathrm{t}\mathbf{)}`$:
$$_t\rho _i=D_i[\rho _i+\beta \rho _i(1\rho _i\rho _j)(U_{ii}+U_{ij}],i,j\{A,B\},ji.$$
(17)
We note here that there is no summation over the same indices, and we also note that in the case of identical species, i.e., $`D_A=D_B`$ and $`U_{AA}=U_{BB}=U_{AB}`$, by adding Eq. (17) for $`i=A`$ and the corresponding one for $`i=B`$ one finds that as expected the density $`\rho =\rho _A+\rho _B`$ satisfies Eq. (5). Assuming that the long-ranged parts of the pair-interactions $`AA`$, $`BB`$, and $`AB`$ are of the same Lennard-Jones form (see Sec. 2) only with different strengths $`U_0^{(ij)}`$, $`i,j\{A,B\}`$, one can repeat the same argument as that following Eq. (5) to reduce Eq. (17) to a local one,
$$_t\rho _i=D_i[\rho _ig\rho _i(1\rho _i\rho _j)(W_{ii}\rho _i+W_{ij}\rho _j],i,j\{A,B\},ji,$$
(18)
where the notation $`W_{ij}=\beta U_0^{(ij)}`$ has been introduced.
We consider a T-junction patterned onto a planar, rectangular substrate of size $`L_x\times 2L_y`$ as described below \[see also Fig. 3(a)\].
The spreading monolayers are extracted from reservoirs of particles $`A`$ and $`B`$, respectively, which maintain constant line densities $`C_0^{A,B}`$ at the lines $`𝒞_1=(y=L_y,0xd)`$ and $`𝒞_2=(y=L_y,0xd)`$, respectively. The domains $`𝒟_1=\{(x,y)|y<h/2x>d\}`$ and $`𝒟_2=\{(x,y)|y>h/2x>d\}`$ represent sites with very low affinity for the fluid particles of either type, such that similarly to the situation in Sec. 3 the boundaries of these domains effectively act as hard walls, confining the spreading onto the two lanes forming the inverted T-junction. Finally, we assume that at the foot $`x=L_x,|y|h/2`$ of the T-junction there is a sink for particles of both species. Under these assumptions, the functions $`\rho _{A,B}(x,y,t)`$ as solutions of Eq. (18) fulfill the initial condition
$`\rho _A(x,y,0)=C_0^A,`$ $`(x,y)𝒞_1,`$
$`\rho _B(x,y,0)=C_0^B,`$ $`(x,y)𝒞_2,`$ (19)
$`\rho _{A,B}(x,y,0)=0,`$ $`\mathrm{otherwise},`$
and the boundary conditions
$`\rho _A(x,y,t)|{}_{𝒞_1}{}^{}=C_0^A,`$
$`\rho _B(x,y,t)|{}_{𝒞_2}{}^{}=C_0^B,`$ (20)
$`\rho _A(L_x,y,t)=\rho _B(L_x,y,t)=0,`$
$`\mathrm{j}_\mathrm{n}^{\mathrm{A}\mathbf{,}\mathrm{B}}\mathbf{|}{}_{\mathbf{}𝓓_{1\mathbf{,}2}}{}^{}\mathbf{=}0\mathbf{,}`$
where the current $`\mathrm{j}^\mathrm{A}`$ ($`\mathrm{j}^\mathrm{B}`$ is obtained by exchanging the labels $`AB`$) is now given by \[see Eq. (18)\]
$$\mathrm{j}^\mathrm{A}\mathbf{=}\mathbf{}𝝆_\mathrm{A}\mathbf{}\mathrm{g}𝝆_\mathrm{A}\mathbf{(}1\mathbf{}𝝆_\mathrm{A}\mathbf{}𝝆_\mathrm{B}\mathbf{)}\mathbf{(}\mathrm{W}_{\mathrm{A}\mathrm{A}}\mathbf{}𝝆_\mathrm{A}\mathbf{+}\mathrm{W}_{\mathrm{A}\mathrm{B}}\mathbf{}𝝆_\mathrm{B}\mathbf{)}\mathbf{.}$$
(21)
In the following we focus on the effect of the A-B interaction on the dynamics of mixing of otherwise identical monolayers, i.e., we choose $`W_{AA}=W_{BB}`$, $`D_A=D_B`$, and $`C_0^A=C_0^B`$; the results discussed in the following correspond to the particular choice $`C_0^A=C_0^B=1`$, while the parameter $`D_A`$ is absorbed into the variable $`\tau =D_At`$. The geometrical parameters are fixed to $`L_x=500`$, $`L_y=20`$, and $`d=h=20`$. The mixing will be characterized by the ratio
$$c_A(\mathrm{r}\mathbf{,}𝝉\mathbf{)}\mathbf{=}\frac{𝝆_\mathrm{A}\mathbf{(}\mathrm{r}\mathbf{,}𝝉\mathbf{)}}{𝝆_\mathrm{A}\mathbf{(}\mathrm{r}\mathbf{,}𝝉\mathbf{)}\mathbf{+}𝝆_\mathrm{B}\mathbf{(}\mathrm{r}\mathbf{,}𝝉\mathbf{)}}$$
(22)
(with the convention $`c_A=0`$ if $`\rho _A+\rho _B=0`$), which is close to 1 in A-rich regions, close to zero in B-rich regions, and close to 1/2 in regions where mixing is accomplished (i.e., $`\rho _A\rho _B0`$).
In Figs. 3(b) and (c) we present typical results for $`c_A(\mathrm{r}\mathbf{,}𝝉\mathbf{)}`$ from numerical integration of the coupled set of equations given by Eq. (18) with the initial and boundary conditions Eqs. (4)-(4) for the case of attractive interactions $`W_{AA}=W_{BB}=0.7`$ and attractive inter-species interaction $`W_{AB}=0.6`$ (b), respectively repulsive inter-species interaction $`W_{AB}=0.6`$ (c). From these figures it is clear that within the stripe forming the leg of the T-junction there is almost perfect mixing.
Surprisingly, at first glance the result seems to be almost independent of the sign of the A-B interaction, and there is only a weak dependence on the strength $`W_{AB}`$ of this interaction, except for the extension of the A-rich and B-rich regions near the corners of the T-junction. However, this behavior can be easily rationalized in view of the fact that, as shown in Ref. , the structure and dynamics of the spreading of a one-component monolayer of particles with inter-particle attraction $`W_{AA}=0.7`$ is well described also by the ”effective boundary force” theory of Burlatsky et al , which disregards interactions within the bulk of the monolayer. The reason for this is that the density in the spreading monolayer is relatively low, except for the region near the reservoir, and thus the system is too dilute to be influenced by the inter-particle attraction. Therefore, in the initial stages of spreading and mixing, the stripe is invaded by low density phases of $`A`$ and $`B`$ which mix independently of their mutual interaction. This scenario is well supported by the comparative analysis of the corresponding density profiles $`\rho _A(\mathrm{r}\mathbf{,}𝝉\mathbf{)}`$ in the cases $`W_{AB}=0.6`$ and $`W_{AB}=0.6`$, respectively, shown in Figs. 3(d) and (e): almost everywhere in the stripe the density of $`A`$ particles is low, in the range of 0 to 0.25, and thus the mixture behaves as a dilute, non-interacting two-dimensional gas. Finally, we note a weak dependence of the extension of the A-rich and B-rich regions near the corners of the T-junction on the sign of the $`W_{AB}`$ interaction \[compare Figs. 3 (b) and (c), respectively (d) and (e)\].
The above scenario holds for all values $`W_{AA}=W_{BB}0.9`$ and $`|W_{AB}|W_{AA}`$ that we have tested, and thus the (tentative) conclusion is that the symmetric T-geometry would ensure practically perfect mixing (but without any possibility to control the spatial extent or the spatial distribution of mixing) for two-component monolayers with a similarity of the interactions between like species. However, further calculations should be carried out before definite conclusions can be drawn concering this issue. For example, for the case $`W_{AA},W_{BB}>1.1`$ we expect sharp interfaces to emerge in each of the two spreading monolayers, which might eliminate the ”mixing through the low density front” mechanism discussed above. Simulations for these ranges of values for the inter-particle interactions turned out to be extremely time consuming, and work is still in progress to elucidate this point.
## 5 Summary
A lattice gas model of interacting particles and the corresponding nonlinear diffusion equation derived from its microscopic dynamics in the continuum limit provide a simple but realistic description of fluid monolayer spreading on a homogeneous substrate. Based on previous results for spreading on a homogeneous substrate, here we have extended this model to address two more complex problems: the spreading of monolayers around obstacles (Sec. 3) and the mixing of monolayers (Sec. 4). These are simple examples of spreading on chemically designed substrates.
For the case of monolayer spreading in the presence of a mesoscopic obstacle, the results obtained from the numerical integration of the nonlinear diffusion equation (Eq. (10)) with initial and boundary conditions given by Eqs. (14)-(3) show that the iso-density lines are bent and pinned by the obstacle during the spreading of the monolayer around it. For a fixed geometry, a fixed density of the reservoir, and fixed substrate-fluid and inter-particle interactions the spatial and temporal extent of this bending in front of and behind the obstacle can be used as measures for the relaxation of the density profile upon passing around the obstacle.
As an example for the mixing of two species in the course of spreading of two monolayers at the merger of two chemical lanes, we have discussed the case of a T-junction geometry. We have focused on the effect of the A-B inter-species interaction on the dynamics of mixing of otherwise identical monolayers. Surprisingly, so far our results lead to the conclusion that the symmetric T-geometry together with the similarity of the same-species interaction ensures practically perfect mixing (but without any possibility to control the spatial extent or the spatial distribution of mixing) for two monolayers of different species, independently of the sign or the strength of the A-B interaction. Only the extension of the A-rich and B-rich regions near the corners of the T-junction exhibits differences. This behavior reflects the fact that in the initial stages of spreading and mixing the stripe is invaded by low density phases of $`A`$ and $`B`$, which mix independently of the inter-species interaction; the resulting mixture is a dilute, quasi non-interacting two-dimensional gas. Because this complete mixing has occurred at early stages, when repulsion does not play a role, and since the continuum equation does not contain any noise terms, demixing or segregation is not observed in the present calculations, although it is expected to occur for strongly repulsive A-B interactions.
## References
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# Analytic Gravitational-Force Calculations for Models of the Kuiper Belt, with Application to the Pioneer Anomaly
## I Introduction
There has long been interest in the gravitational force that could be produced by the Kuiper Belt boss . It has been observed that total masses of much more than an Earth mass, $`M_{}`$, would lead to conflicts with orbital observations. (See, e.g., Refs. boss ; KBjda and Sec. VII-E of pioprd .) Further, it has also been calculated that a Kuiper-Belt ring with a mass of this magnitude could not explain an acceleration the size of the Pioneer anomaly boss ; KBjda ; pioprd . This anomaly pioprd ; pioprl is the apparent unmodeled constant acceleration of the Pioneer spacecraft, observed between $`2070`$ Astronomical Units (AU), of magnitude
$$a_P(20\mathrm{A}\mathrm{U}<r<70\mathrm{A}\mathrm{U})=(8.74\pm 1.33)\times 10^8\mathrm{cm}/\mathrm{s}^2$$
(1)
which is directed approximately towards the Sun.
Even so, this type of Kuiper-Belt mechanism has remained a fascinating one as a possible explanation of the anomaly. In particular, it has recently been proposed diego that gravitation from the Kuiper Belt, modeled by a cylindrically-symmetric ring of matter whose density goes as
$$\rho _1(p)=\frac{\rho _1L}{p},p=\sqrt{x^2+y^2},$$
(2)
where
$$\rho _1=1.74\times 10^{16}\mathrm{g}/\mathrm{cm}^3,L=20\mathrm{AU},$$
(3)
can explain the constant anomaly. The ring has a width
$$R_1=20\mathrm{AU}p100\mathrm{AU}=R_2$$
(4)
and a thickness
$$2D=2\mathrm{AU}.$$
(5)
The mass is thus
$$_{ring}=4\pi \rho _1LD(R_2R_1)=1.17\times 10^{28}\mathrm{g}=1.96M_{}.$$
(6)
This proposal is somewhat surprising, given the observations noted above. However, one is thereby motivated to take a different looks at the problem orfeu . Here we do so emphasizing analytic calculations. This will help to better understand the underlying physics of the situation.
To start, although it is well-known that a spherically symmetric ball with a density that goes as $`1/r`$ can produce a constant acceleration within the ball, there only is a constant acceleration from a complete spherical ball, not from a shell. Therefore, as we emphasize in the next section, with only a cylinder ring, not even a cylindrical disk, satisfying a constant acceleration is doubly hard to do. Specifically, it can not come from an exact cylindrically-symmetric $`1/p`$ density. Indeed, although the appeal to Gauss’ Law in Eq. (3) of diego is correct, the argument that Eq. (4) of diego implies there will be a constant acceleration within the ring is not exact. We will demonstrate this by specific analytical calculation.
Before continuing, we note again that the mass of the model belt of Ref. diego appears to be somewhat high, as has been determined elsewhere boss ; KBjda ; pioprd . Further, it is known that the amount of dust is much smaller than this, and the gravitational mass of the Kuiper Belt is dominated by large rocks and ices. The interplanetary dust is actually supplied by collisions between the rocks and ices and lives for only of order 100,000 years in the inner solar system, an order of magnitude longer in the outer solar system. Further, the dominant mass of the rocks and ices is overwhelming subject to gravity and not other forces. Hence, there tend to be resonant concentrations in it vs. a smooth distribution mann1 -drag .
In Section II we will describe the gravity of spherical balls and shells. This is followed by an introduction to the gravity of cylindrically symmetric disks and rings in Section III. (These objects are examined in both the complete 3-dimensional framework and also in the “thin-ring” approximation, where the distribution in the $`z`$ direction is a $`\delta `$-function.) In Section IV we apply the “thin-ring” approximation to both the $`1/p`$ model and the Boss-Peale model boss . We then go on to full 3-dimensional calculations. In Sections V, VI, and VII we discuss, respectively, the $`1/p`$-density cylindrical ring, the constant-density cylindrical ring, and the $`1/r^2`$-density wedge (as well as the $`1/p^2`$-density “thin ring”). We end with a discussion where we compare the results. In particular, we compare the accelerations produced by the 3-dimensional $`1/p`$-, $`1/r^2`$-, and constant-density rings, as well as those from the Boss-Peale and $`1/p^2`$ “thin-rings.”
We find, as expected, that neither the magnitude nor the shape of the Pioneer anomaly can be reproduced. (For comparison, in our numerical plots we will adhere as much as possible to the model parameters of Eq. (3). However, since the basic formulae are analytic, they can be renormalized at will.)
## II Spherical balls and shells
The $`1/r^2`$ gravitational force law yields that any spherically symmetric distribution with total mass $``$ exerts a force outside that distribution that is proportional to the total mass divided by the square of the distance to the center of symmetry: $`G/r^2`$. Contrarily, if the observation point is inside a spherical distribution of mass, no force is exerted.
This is an important result for understanding the effects of a general spherically symmetric density distribution, $`\rho (r)`$. Since we are heading towards the $`1/r`$ distribution, consider density distributions that go as
$$\rho (r)\frac{\rho _n(r)L^n}{r^n},\mathrm{}n\mathrm{}.$$
(7)
Here $`\{\rho _n,L\}`$ give the overall normalizations in terms of some density and length scale. These types of densities have long been studied by geophysicists. They often like to think in terms of spherical distributions and shells of the Earth having different functional dependences and thus causing different gravity signals dzi81 ; geo . But note: We are talking about spherical shells, not cylindrical rings.
In the present study, we will concentrate on the distributions for $`n=\{0,1,2\}`$, the constant, $`1/r`$, and $`1/r^2`$ distributions. Specifically, start with the $`1/r`$ distribution,
$$\rho _1(r)=\frac{\rho _1L}{r}.$$
(8)
It has a total mass out to a radius $`R`$ of
$$_{ball}(R)=2\pi \rho _1LR^2.$$
(9)
(Of course, if the density distribution went to infinity there would be infinite mass.) From the spherical symmetry condition mentioned before, we have that interior and exterior to the sphere
$`a_{ball}(r<R)`$ $`=`$ $`{\displaystyle \frac{G_{ball}(r)}{r^2}}=G2\pi \rho _1L,`$ (10)
$`a_{ball}(r>R)`$ $`=`$ $`{\displaystyle \frac{G_{ball}(R)}{r^2}}={\displaystyle \frac{G2\pi \rho _1LR^2}{r^2}}.`$ (11)
That is, there is a constant acceleration inside the ball and the ordinary Newtonian inverse-square force outside the ball. Even so, there remains a singularity at the origin since there the acceleration is a non-zero constant pointing radially in from all directions.
If we now use the parameters of Ref. diego given in Eq. (3) above, $`\rho _1=1.74\times 10^{16}`$ g/cm<sup>3</sup> and $`L=20`$ AU, then even the spherical ball of Eq. (10) would only produce an acceleration of magnitude
$`a_{ball}(r<R)`$ $`=`$ $`C_{ball}=(2\pi G\rho _1L)`$ (12)
$`=`$ $`2.18\times 10^8\mathrm{cm}/\mathrm{s}^2.`$
But this is smaller than $`a_P`$! So, if an entire ball of this density can not cause the Pioneer anomaly, how can a disk, let alone a ring?
To continue, what if this were only a spherical shell (from $`R_1=20`$ AU to $`R_2=100`$ AU)? Then, even inside the shell the acceleration would not be constant. By subtracting out the gravitational attraction of the mass interior to radius $`R_1`$ the acceleration is
$`a_{shell}(0<r<R_1)`$ $`=`$ $`0,`$ (13)
$`a_{shell}(R_1<r<R_2)`$ $`=`$ $`G2\pi \rho _1L+{\displaystyle \frac{G_{ball}(R_1)}{r^2}}`$ (14)
$`=`$ $`G2\pi \rho _1L+{\displaystyle \frac{G2\pi \rho _1LR_1^2}{r^2}},`$
$`a_{shell}(r>R_2)`$ $`=`$ $`{\displaystyle \frac{G2\pi \rho _1L(R_2^2R_1^2)}{r^2}},`$ (15)
where we write
$`a_{shell}(r)`$ $``$ $`(2\pi G\rho _1L)g_{shell}(r)=C_{ball}g_{shell}(r),`$ (16)
$`_{shell}`$ $`=`$ $`2\pi \rho _1L(R_2^2R_1^2)=60_{ring}`$ (17)
$`=`$ $`7.03\times 10^{29}\mathrm{g}.`$
Therefore, there is a constant acceleration towards the center of a spherical $`1/r`$-density distribution of matter given by Eq. (8) only if the mass distribution goes all the way into the origin; that is, if it is a spherical ball, not a spherical shell. In Figure 1 we show $`a_{shell}(r)`$ vs $`r`$ for the values $`\{R_1,R_2\}=\{20,100\}`$ AU. This figure will be useful for comparison when we go to rings.
Particular values of the acceleration are
$$10^8a_{shell}(\{10,60,120\}\mathrm{AU})=\{0,1.94,1.45\}\mathrm{cm}/\mathrm{s}^2.$$
(18)
However, even here with only the first 20 AU of the 100 AU ball deleted, the acceleration varies by an order 10% in the outer half of the shell and rapidly decreases to zero interior to that.
## III Cylindrical disks and rings
### III.1 Full 3-d disks and rings
Now we go on to disks and rings. We use a method inspired by techniques to analyze cylinder cylindrically-symmetric objects in laboratory big-G experiments heyl -newman . The general potential functional and acceleration from a cylindrical symmetric ring are
$`𝒱(r)`$ $`=`$ $`V(r)/m_{test},`$ (19)
$`𝒱(r)`$ $`=`$ $`G{\displaystyle _D^{+D}}dz{\displaystyle _{R_1}^{R_2}}dpp\rho (p)\times `$ (20)
$`{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\varphi }{\sqrt{p^2+r_x^22r_xp\mathrm{cos}\varphi +(r_zz)^2}}},`$
$`a_x(r)`$ $`=`$ $`G{\displaystyle \frac{d}{dr_x}}\left[𝒱(r)\right],`$ (21)
$`a_z(r)`$ $`=`$ $`G{\displaystyle \frac{d}{dr_z}}\left[𝒱(r)\right].`$ (22)
In the above, by convention we take the component of the direction to the test mass in the plane of the ecliptic to be along the $`x`$ axis: $`r_{\{x,y\}}r_x`$. This is useful since we will concentrate on the case of axial symmetry. We also observe that the $`z`$-component of the acceleration in Eq. (22), for general positions out of the ecliptic, is easier to handle lass ; BT in the “thin-ring” approximation of the next subsection.
We denote these various choices by:
$$𝐫(r_x,0,r_z)_{ecliptic}(r,0,0).$$
(23)
(Note for future reference that, with cylindrical symmetry, the volume element, $`p`$, cancels the $`(1/p)`$ of a $`\rho _1(p)`$ density function.)
### III.2 “Thin-ring” approximation
#### III.2.1 General thin rings
As an initial step, we start in the next section by using an analytic approximation,
$$\rho (r)2D\delta (z)\rho (p).$$
(24)
We can do this because $`z`$ is generally small compared to $`p`$ so the change in the overall result should be small and still symmetric about the z axis.
This yields
$`𝒱_{thin}(r)`$ $`=`$ $`2GD{\displaystyle _{R_1}^{R_2}}dpp\rho (p)\times `$ (25)
$`{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\varphi }{\sqrt{p^2+r_x^22r_xp\mathrm{cos}\varphi +r_z^2}}},`$
$`a_{thin}(r)`$ $`=`$ $`2GD{\displaystyle \frac{d}{dr_x}}[{\displaystyle _{R_1}^{R_2}}dpp\rho (p)\times `$ (26)
$`{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\varphi }{\sqrt{p^2+r_x^22r_xp\mathrm{cos}\varphi +r_z^2}}}].`$
#### III.2.2 Taking the $`𝐫`$-derivative first
One tack that can be taken (and will be in Sections IV.1 and IV.2 below) is to first perform the $`r_x`$-derivative in Eq. (26),
$`a_{thin}(r)`$ $`=`$ $`4GD{\displaystyle _{R_1}^{R_2}}dpp\rho (p)\times `$ (27)
$`{\displaystyle _0^\pi }{\displaystyle \frac{d\varphi (r_xp\mathrm{cos}\varphi )}{[p^2+r_x^22r_xp\mathrm{cos}\varphi +r_z^2]^{3/2}}},`$
and then do the $`\varphi `$-integral. Going to the plane of the ecliptic, $`r_z0`$, the result is
$`a_{thin}(r)`$ $`=`$ $`4GD{\displaystyle _{R_1}^{R_2}}𝑑pp\rho (p)\left[{\displaystyle \frac{𝐊\left(\sqrt{\frac{4pr}{r^22pr+p^2}}\right)}{r\sqrt{r^22pr+p^2}}}+{\displaystyle \frac{(rp)𝐄\left(\sqrt{\frac{4pr}{r^22pr+p^2}}\right)}{r(r+p)\sqrt{r^22pr+p^2}}}\right]`$ (28)
$`=`$ $`4GD{\displaystyle _{R_1}^{R_2}}𝑑pp\rho (p)\left[{\displaystyle \frac{𝐊\left(\sqrt{\frac{4pr}{r^2+2pr+p^2}}\right)}{r(r+p)}}+{\displaystyle \frac{𝐄\left(\sqrt{\frac{4pr}{r^2+2pr+p^2}}\right)}{r(rp)}}\right],`$ (29)
where the last two equalities are related by 8.127 of Ref. GR and the complete elliptic integrals of the first and second kind (see 8.113 and 8.114 in GR ) are
$`𝐊(t)`$ $``$ $`K(t^2)={\displaystyle \frac{\pi }{2}}F({\displaystyle \frac{1}{2}},{\displaystyle \frac{1}{2}};1;t^2)={\displaystyle \frac{\pi }{2}}\left(1+{\displaystyle \frac{t^2}{4}}+{\displaystyle \frac{9t^4}{64}}+\mathrm{}+\left[{\displaystyle \frac{(2n1)!!}{2^nn!}}\right]^2t^{2n}+\mathrm{}\right),`$ (30)
$`𝐄(t)`$ $``$ $`E(t^2)={\displaystyle \frac{\pi }{2}}F({\displaystyle \frac{1}{2}},{\displaystyle \frac{1}{2}};1;t^2)={\displaystyle \frac{\pi }{2}}\left(1{\displaystyle \frac{t^2}{4}}{\displaystyle \frac{3t^4}{64}}\mathrm{}\left[{\displaystyle \frac{(2n1)!!}{2^nn!}}\right]{\displaystyle \frac{t^{2n}}{2n1}}\mathrm{}\right).`$ (31)
This yields a physically intuitive $`p`$-integration that can be handled numerically boss . We will use Eq. (29) in Sections IV.2 and IV.3 below.
## IV Specific thin rings
### IV.1 Analytic, thin-ring, $`\mathrm{𝟏}/𝐩`$-density model
Returning to Section III.2.1, it turns out that the thin-ring, $`1/p`$-density problem is analytically solvable. If one does the $`\varphi `$-integral before the r-differentiation in Eq. (26) one can also do the second integral. (Again note, for this $`1/p`$-density case, the density function cancels the $`p`$ in the volume element, making the integrals simpler.) Proceeding, the potential functional in the plane of the ecliptic is
$$𝒱_{T/p}(r)=G\rho _12DL_{R_1}^{R_2}𝑑p_0^{2\pi }\frac{d\varphi }{\sqrt{p^2+r^22rp\mathrm{cos}\varphi }}.$$
(32)
The $`\varphi `$ integral is analytic and is (3.674.1 in GR )
$`I_\varphi (r>p)`$ $`=`$ $`{\displaystyle \frac{4}{r}}𝐊\left({\displaystyle \frac{p}{r}}\right),`$ (33)
$`I_\varphi (r<p)`$ $`=`$ $`{\displaystyle \frac{4}{p}}𝐊\left({\displaystyle \frac{r}{p}}\right).`$ (34)
(Eqs. (32) and (33) demonstrate that for very large $`r`$ the potential goes to $`GM_{ring}/r`$, as it should.)
This means that the potentials outside, within, and inside of the ring are
$`𝒱_{T/p}(R_2<r)=8G\rho _1LD{\displaystyle _{R_1}^{R_2}}{\displaystyle \frac{dp}{r}}𝐊\left({\displaystyle \frac{p}{r}}\right),`$
$`𝒱_{T/p}(R_1<r<R_2)=8G\rho _1LD[{\displaystyle _{R_1}^r}{\displaystyle \frac{dp}{r}}𝐊\left({\displaystyle \frac{p}{r}}\right)`$
$`+{\displaystyle _r^{R_2}}{\displaystyle \frac{dp}{p}}𝐊\left({\displaystyle \frac{r}{p}}\right)],`$ (36)
$`𝒱_{T/p}(r<R_1)=8G\rho _1LD{\displaystyle _{R_1}^{R_2}}{\displaystyle \frac{dp}{p}}𝐊\left({\displaystyle \frac{r}{p}}\right).`$ (37)
Changing variables to $`t=p/r`$ or $`r/p`$, respectively, and using the properties of the complete elliptic integral, the acceleration ($`a_{T/p}=d𝒱_{T/p}/dr`$) is
$`a_{T/p}(r)=C_1g_{T/p}(r),`$ (38)
$`C_1=8G\rho _1L=(4/\pi )C_{ball}=2.779\times 10^8\mathrm{cm}/\mathrm{s}^2,`$ (39)
$`g_{T/p}(R_2<r)={\displaystyle \frac{DR_2}{r^2}}𝐊\left({\displaystyle \frac{R_2}{r}}\right){\displaystyle \frac{DR_1}{r^2}}𝐊\left({\displaystyle \frac{R_1}{r}}\right),`$ (40)
$`g_{T/p}(R_1<r<R_2)={\displaystyle \frac{D}{r}}𝐊\left({\displaystyle \frac{r}{R_2}}\right){\displaystyle \frac{DR_1}{r^2}}𝐊\left({\displaystyle \frac{R_1}{r}}\right),`$
(41)
$`g_{T/p}(r<R_1)={\displaystyle \frac{D}{r}}𝐊\left({\displaystyle \frac{r}{R_2}}\right){\displaystyle \frac{D}{r}}𝐊\left({\displaystyle \frac{r}{R_1}}\right).`$ (42)
This acceleration is not a constant for $`(R_1<r<R_2)`$.
Putting the remaining distances in terms of AU, in Figure 2 we plot $`a_{T/p}(r)`$ vs. $`r`$ using the parameters of Ref. diego . One can note the general features. Most importantly, the size of the acceleration within this model of the Kuiper Belt is about a factor of 100 smaller than the anomaly! In particular, specific values of the acceleration are
$`a_{T/p}(\{10,60,120\}\mathrm{AU})=`$ (43)
$`\{0.0309,+0.0610,+0.0338\}\times 10^8\mathrm{cm}/\mathrm{s}^2,`$
which can be compared to the values from a shell given in Eq. (18). The acceleration within the ring is of order 40 times smaller than that within the shell.
Observe that $`a_{T/p}(r)`$ manifestly has other appropriate physical properties. First, $`a_{T/p}(R_2r)G_{ring}/r^2`$. Next, as it should on physical grounds, $`a_{T/p}(r0)0_{}`$. Analytically, Eqs. (30), (38), and (42) show that $`a_{T/p}(r)`$ is slightly negative as $`r0`$ and goes to zero in the limit.
One also sees the breakdowns at $`r=\{R_2,R_1\}`$ where the $`𝐅`$ are singular because the arguments are unity. (Here and later we will cut off the heights of the 2-d spikes.) As we will see, these singularities result from having only a 2-d approximation for the non-smooth (hard-edged) ring. When the density is continuous in the $`p`$ variable the spike singularity in the acceleration disappears, even for 2-d problems. When the problem is 3-d, the spikes become finite cusps. (See Section V.)
As observed, far out $`a_{T/p}(r)`$ goes as $`1/r^2`$. As one comes in, approaches, and then passes $`r=R_2`$, the quantity $`a_{T/p}(r)`$ starts to decrease since less mass is interior to the test point. Within the interior of the ring, for a short distance $`a_{T/p}(r)`$ is “roughly,” but not exactly, flat. (It will be less constant in the true 3-d calculation.) Further, as one gets closer to $`R_1`$ the acceleration changes sign because more mass begins pulling out rather than in. As predicted one sees that $`a_{T/p}(r)`$ is slightly negative as $`r0`$ and it goes to zero at the origin.
### IV.2 Another thin-ring, $`\mathrm{𝟏}/𝐩`$-density calculation
We demonstrate here that an equivalent result for the $`1/p`$-density can be obtained by the method of Section III.2.2. This demonstration illuminates this method which will be useful in the following subsection.
If the $`1/p`$-density given in Eq. (2) is placed in Eq. (29), this yields the acceleration (again $`D`$ will be $`1`$ AU)
$`a_{BP/p}(r)=[(4GL)\rho _1]D{\displaystyle _{20}^{100}}dp\times `$ (44)
$`\left[{\displaystyle \frac{𝐊\left(\sqrt{\frac{4pr}{r^2+2pr+p^2}}\right)}{r(r+p)}}+{\displaystyle \frac{𝐄\left(\sqrt{\frac{4pr}{r^2+2pr+p^2}}\right)}{r(rp)}}\right]`$
$``$ $`[(8GL\rho _1)/2]g_{BP/p}(r)`$ (45)
$`=`$ $`(C_1/2)g_{BP/p}(r)=(2/\pi )C_{ball}g_{BP/p}(r).`$ (46)
The numerical integration yielding $`g_{BP/p}(r)`$ has to deal with integrable singularities at $`p=r`$, which exist because there the argument of $`𝐊`$ is unity. By avoiding the singularities, the integral is doable, except for the two singularities coming from the discontinuous nature of the ring’s density at the boundaries. The result agrees numerically with the result in the previous subsection. That is,
$$a_{BP/p}(r)=a_{T/p}(r),g_{BP/p}(r)=2g_{T/p}(r).$$
(47)
### IV.3 The Boss-Peale model
Eq. (29) is the integral used by Boss and Peale boss to study gravity from a smooth cylindrical mass distribution of the form
$`\rho _{BP}(p)={\displaystyle \frac{\rho _0^{BP}(pA)^2}{D^2}}\mathrm{exp}\left[{\displaystyle \frac{(pA)}{5}}\right],`$ (48)
$`A=50\mathrm{AU}p100\mathrm{AU}=B,D=1\mathrm{AU}.`$ (49)
For comparison we take this model to have the same mass, $`_{ring}`$, given in Eq. (6). Therefore,
$`_{ring}`$ $`=`$ $`4\pi D\rho _0^{BP}D^2{\displaystyle _{50}^{100}}dpp(p50)^2\times `$ (50)
$`\mathrm{exp}\left[{\displaystyle \frac{(pA)}{5}}\right]`$
$`=`$ $`4\pi D\rho _0^{BP}D^25^4[25.8826].`$
(If one makes the approximation that the upper limit of the integral goes to infinity, then the last term in the second line would be $`26=[\mathrm{\Gamma }(4)+10\mathrm{\Gamma }(3)]`$.) Therefore,
$$\rho _0^{BP}=\frac{64}{(25.8826)25}\rho _1=(0.172)\times 10^{16}\mathrm{g}/\mathrm{cm}^3.$$
(51)
If we place this density in Eq. (29) we obtain
$`a_{BP}(r)=C_{BP}g_{BP}(r),`$ (52)
$`C_{BP}`$ $`=`$ $`(4GD\rho _0^{BP})={\displaystyle \frac{8}{5^3(25.8826)}}C_1`$ (53)
$`=`$ $`0.002473C_1=(0.00687)\times 10^8\mathrm{cm}/\mathrm{s}^2,`$
where the quantity $`g_{BP}(r)`$ is
$`g_{BP}(r)={\displaystyle _{50}^{100}}dpp(p50)^2\mathrm{exp}[{\displaystyle \frac{(p50)}{5}}]\times `$
$`\left[{\displaystyle \frac{𝐊\left(\sqrt{\frac{4pr}{r^2+2pr+p^2}}\right)}{r(r+p)}}+{\displaystyle \frac{𝐄\left(\sqrt{\frac{4pr}{r^2+2pr+p^2}}\right)}{r(rp)}}\right].`$ (54)
As in the last subsection, $`g_{BP}(r)`$ can be integrated numerically boss , but with difficulty because of the integrable singularities when $`r=p`$. The result for $`a_{BP}(r)`$ is shown in Figure 3, which agrees with Figure 1 of Ref. boss (except for the small, narrow spike at $`B=100`$ – see below).
Particular values of the acceleration are
$``$ $`a_{BP}(\{10,53,73,120\}\mathrm{AU})=`$
$`\{0.00686,0.212,+0.159,+0.0325\}10^8\mathrm{cm}/\mathrm{s}^2.`$ (55)
These values, and the shape of Figure 3 reflect the different type of density profile of this ring. Note that the curve for $`a_{BP}(r)`$ is smooth when $`r=50`$. This is because the density varies continuously from zero at this point. On the other hand, note the small, narrow spike at $`r=100`$, which occurs since the ”thin” ring abruptly ends there with the density $`\rho _{BP}(p)/\rho _{BP}^0`$ going discontinuously from $`(2500\mathrm{exp}[10])=0.113`$ to zero. If the ring density is allowed to smoothly continue on past $`r=100`$, decreasing exponentially out to infinity, the spike disappears and the resulting $`a_{BP_{\mathrm{}}}(r)`$ becomes very slightly higher (lower) in magnitude than $`a_{BP}(r)`$ going somewhat further out (in) from the position of the spike.
A comparison of the normalized acceleration, $`a_{BP}(r)`$, with that for other models will be given in Section VIII
## V 3-d, cylindrical-coordinate, $`(\mathrm{𝟏}/𝐩)`$-density ring
Now we calculate the acceleration from the $`(1/p)`$-density in the 3-d case. Begin with the complete, exact, 3-dimensional integral defined in Eqs. (20) and (21) with the ring $`(1/p)`$-density of Eq. (2):
$$a_{1/p}(r)=G\left[\frac{d}{dr_x}\right]_0^{2\pi }𝑑\varphi _D^{+D}𝑑z_{R_1}^{R_2}\frac{dpp\rho _1L}{p\sqrt{p^2+r_x^22rp\mathrm{cos}\varphi +(zr_z)^2}}.$$
(56)
Going to the plane of the ecliptic, performing the $`p`$-integration (which is easy since the density cancels the volume element), and then doing the $`r`$-derivative yields
$`a_{1/p}(r)`$ $``$ $`(C_1/4)g_{1/p}(r)`$ (57)
$`a_{1/p}(r)`$ $`=`$ $`{\displaystyle \frac{C_1}{4}}\left[{\displaystyle \frac{d}{dr}}\right]{\displaystyle _0^\pi }d\varphi {\displaystyle _D^{+D}}dz\times \mathrm{ln}[pr\mathrm{cos}\varphi +\sqrt{p^2+r^2+z^22pr\mathrm{cos}\varphi }]_{R_1}^{R_2}`$ (58)
$`=`$ $`{\displaystyle \frac{C_1}{4}}{\displaystyle _0^\pi }𝑑\varphi {\displaystyle _D^{+D}}𝑑z[\mathrm{\Phi }(r,R_2,z,\varphi )\mathrm{\Phi }(r,R_1,z,\varphi )],`$
$`\mathrm{\Phi }(r,R,z,\varphi )`$ $`=`$ $`{\displaystyle \frac{\mathrm{cos}\varphi +(R\mathrm{cos}\varphi r)/S}{(Rr\mathrm{cos}\varphi )+S}}`$ (59)
$`=`$ $`\left[{\displaystyle \frac{r\mathrm{sin}^2\varphi }{z^2+r^2\mathrm{sin}^2\varphi }}\right]+\left[{\displaystyle \frac{S\mathrm{cos}\varphi +[(R^2+r^2)\mathrm{cos}\varphi +pR(1+\mathrm{cos}^2\varphi )]/S}{z^2+r^2\mathrm{sin}^2\varphi }}\right],`$
$`S`$ $``$ $`\sqrt{R^2+r^22rR\mathrm{cos}\varphi +z^2}.`$ (60)
The $`z`$-integration can be done analytically using the two sets of square brackets in Eq. (59) separately, with the complicated second piece adding an additional part to the first term. This yields
$`H(r,R,Z,\varphi )`$ $`=`$ $`2\mathrm{sin}\varphi \mathrm{tan}^1\left[{\displaystyle \frac{Z}{r\mathrm{sin}\varphi }}\right]+\mathrm{cos}\varphi \mathrm{ln}[Z+S]+\mathrm{sin}\varphi \mathrm{tan}^1\left[{\displaystyle \frac{rZ\mathrm{sin}\varphi }{R^2+r^22Rr\mathrm{cos}\varphi +(R+r\mathrm{cos}\varphi )S}}\right],`$
$`S`$ $``$ $`\sqrt{R^2+r^22rR\mathrm{cos}\varphi +Z^2}.`$ (62)
Although it is technically possible to do the $`\varphi `$-integration, the result is so complicated that it is preferable to do the final integral numerically. The result,
$$g_{1/p}(r)=_0^\pi 𝑑\varphi \left[H(r,R_2,D,\varphi )H(r,R_2,D,\varphi )H(r,R_1,D,\varphi )+H(r,R_1,D,\varphi )\right],$$
(63)
is used to obtain $`a_{1/p}(r)`$, which is shown in Figure 4. (The numerical singularities to be overcome occur when $`(rR_{\{1,2\}}\mathrm{cos}\varphi )=0`$.
This figure again has the correct behaviour. It is a more delicate version of the “thin ring” shown in Fig. 2. The most noticeable change from the “thin ring” is that the spikes of $`a_{T/p}(r)`$ near $`\{R_1,R_2\}`$ in Figure 2 become finite cusps at $`\{R_1,R_2\}`$ of $`a_{1/p}(r)`$ in Figure 4. The cusps are also less extreme compared to the spikes. This is because, for the 3-d ring, all the near-by mass is not at a point on the ecliptic, but along a line perpendicular to it. The proper limit can be seen by evaluating both the 2-d and 3-d forms as $`r`$ becomes large. By $`r=1000`$ the two forms already agree to three significant figures.
To summarize: A $`1/p`$-density potential in a ring does not produce a constant acceleration within the ring.
A comparison of the normalized acceleration, $`a_{1/p}(r)`$, with that for other models will be given in Section VIII.
## VI Cartesian, constant-density ring
We next consider a constant density disk. This is of interest for both physical and mathematical comparisons. We use cartesian coordinates because for cartesian coordinates the volume element is unity. Therefore, a constant density has the simplest integrals with these coordinates. (We already observed how the $`1/\sqrt{x^2+y^2}`$ density cancels the $`\sqrt{x^2+y^2}`$ volume element in cylindrical coordinates.) This current calculation is similar to that used in Ref. cylinder to study the metrology of a solid cylinder for big-G Cavendish experiments.
To settle on $`\rho _0`$ we take the same total mass and shape as the $`1/p`$ ring. This means
$$\rho _0=\rho _1\frac{2L}{R_1+R_2}=\rho _1/3.$$
(64)
Now proceed by using Eq. (20), giving
$$𝒱_{Con}(r)=G\rho _0\left[_{R_2}^{R_2}𝑑y_{\sqrt{R_2^2y^2}}^{\sqrt{R_2^2y^2}}𝑑x_{R_1}^{R_1}𝑑y_{\sqrt{R_1^2y^2}}^{\sqrt{R_1^2y^2}}𝑑x\right]_D^{+D}\frac{dz}{\sqrt{(xr_x)^2+y^2+(zr_z)^2}}.$$
(65)
The two integrals represent the gravitational effect of a disk of radius $`R_2`$ minus the effect of a disk of radius $`R_1`$, thus yielding a ring. Again, in the plane of the ecliptic ($`r_z=0`$) the acceleration is obtained by taking the negative of the derivative of the integrand with respect to $`r`$:
$$\frac{d}{dr}\frac{1}{[(rx)^2+y^2+z^2]^{1/2}}=\frac{rx}{[(rx)^2+y^2+z^2]^{3/2}}.$$
(66)
(Note that since one is taking the derivative of the square root of a square, one must be careful that the correct over-all sign emerges.)
Now do the integral with respect to $`z`$. This yields
$$I_z=\left[\frac{z(rx)}{[(rx)^2+y^2][z^2+(rx)^2+y^2]^{1/2}}\right]_D^D.$$
(67)
Thus,
$$a_{Con}(r)=G\rho _0\left[_{R_2}^{R_2}𝑑y_{\sqrt{R_2^2y^2}}^{\sqrt{R_2^2y^2}}𝑑x_{R_1}^{R_1}𝑑y_{\sqrt{R_1^2y^2}}^{\sqrt{R_1^2y^2}}𝑑x\right]\frac{2D(rx)}{[(rx)^2+y^2][D^2+(rx)^2+y^2]^{1/2}}.$$
(68)
The $`x`$ integral is
$$I_x=\mathrm{ln}[+D+\sqrt{D^2+(rx)^2+y^2}]_{R_1}^{R_2}\mathrm{ln}[D+\sqrt{D^2+(rx)^2+y^2}]_{R_1}^{R_2},$$
(69)
so
$`a_{Con}(r)`$ $`=`$ $`G\rho _0\left[{\displaystyle _{R_2}^{R_2}}𝑑yF(r,y,R_2,D){\displaystyle _{R_1}^{R_1}}𝑑yF(r,y,R_1,D)\right],`$ (70)
$`F(r,y,R,D)`$ $`=`$ $`\mathrm{ln}\left\{\left[{\displaystyle \frac{[+D+\sqrt{D^2+R^2+r^22r\sqrt{R^2y^2}}]}{[D+\sqrt{D^2+R^2+r^22r\sqrt{R^2y^2}}]}}\right]\left[{\displaystyle \frac{[D+\sqrt{D^2+R^2+r^2+2r\sqrt{R^2y^2}}]}{[+D+\sqrt{D^2+R^2+r^2+2r\sqrt{R^2y^2}}]}}\right]\right\}.`$ (71)
This final integral can be done analytically using involved transformations similar to those used in Ref. cylinder . But the end result is very complicated. Therefore, for clarity, a simple 1-dimensional numerical integral will be used. (As a result we leave unaddressed the implications of the relative sizes of $`r`$ vs. $`\{R_1,R_2\}`$, which implications can play in the analytic form of this final integral.) We change all units to AU, e.g., change the variable $`y`$ to $`t=y/D`$ and multiply the external constants by the same $`D=1`$ AU. Then,
$`a_{Con}`$ $`=`$ $`C_0[{\displaystyle _{100}^{100}}dtF(r,t,100,1)`$ (72)
$`{\displaystyle _{20}^{20}}dtF(r,t,20,1)]`$
$`=`$ $`C_0g_{Con}(r),`$ (73)
$`C_0`$ $`=`$ $`G\rho _0D=C_1/480=0.00579\times 10^8\mathrm{cm}/\mathrm{s}^2.`$ (74)
In Figure 5 we show $`a_{Con}(r)`$. Again we see the correct general behaviour. With the 3-d calculation, the cusps at the discontinuous boundaries of the ring are large, but finite and hence physical. Interesting values of the acceleration are
$`a_{Con}(\{10,20,60,100,120\}\mathrm{AU})\times 10^8=`$
$`\{0.0165,0.0870,+0.03146,+0.130,+0.371\}\mathrm{cm}/\mathrm{s}^2.`$
(75)
Since the total mass is the same as for the $`1/p`$ ring, the acceleration should tend to the same limit as $`r`$ gets large, and it does.
A comparison of the normalized acceleration, $`a_{Con}(r)`$, with that for other models will be given in Section VIII.
## VII Wedge $`\mathrm{𝟏}/𝐫^\mathrm{𝟐}`$ (thin-ring $`\mathrm{𝟏}/𝐩^\mathrm{𝟐}`$) density
### VII.1 Wedge configuration
Now we consider a wedge-shaped slice with the spherical density
$$p_2(r)=\frac{\rho _2L^2}{r^2}.$$
(76)
As before, the slice goes between $`R_1`$ and $`R_2`$, except in spherical distance from the origin. The opening wedge angle is
$$\theta _0=\mathrm{tan}^1(D/R_1)=0.049958\mathrm{radians}.$$
(77)
Keeping the mass of the slice the same,
$`_{ring}`$ $`=`$ $`2\pi (2\delta )\rho _2L^2(R_2R_1),`$ (78)
$`\delta `$ $`=`$ $`\mathrm{sin}\theta _0=1/\sqrt{401}=0.049938,`$ (79)
one has,
$$\rho _2=D/(\delta L)\rho _1\beta \rho _1=(1.0012)\rho _1.$$
(80)
In the plane of the ecliptic the acceleration from the wedge is
$`a_{1/r^2}(r)=G\left[{\displaystyle \frac{d}{dr}}\right]{\displaystyle _{\pi /2\theta _0}^{\pi /2+\theta _0}}d\theta \mathrm{sin}\theta {\displaystyle _0^{2\pi }}d\varphi \times `$
$`{\displaystyle _{R_1}^{R_2}}{\displaystyle \frac{\rho _2L^2}{t^2}}{\displaystyle \frac{t^2dt}{\sqrt{t^2+r^22rt\mathrm{cos}\varphi \mathrm{sin}\theta }}}.`$ (81)
Because the density-functional again cancels the volume element, the $`t`$-integral yields
$$_t=\mathrm{ln}[tr\mathrm{sin}\theta \mathrm{cos}\varphi +\sqrt{t^2+r^22tr\mathrm{sin}\varphi \mathrm{cos}\theta }]_{R_1}^{R_2}.$$
(82)
Now taking the negative of the $`r`$-derivative yields
$`a_{1/r^2}(r)=G\rho _2L^2{\displaystyle _{\pi /2\theta _0}^{\pi /2+\theta _0}}d\theta \times `$
$`\mathrm{sin}\theta {\displaystyle _0^{2\pi }}𝑑\varphi U(r,R_1,R_2,\theta ,\varphi ),`$ (83)
$`U(r,R_1,R_2,\theta ,\varphi )=`$
$`\left[{\displaystyle \frac{(\mathrm{sin}\theta \mathrm{cos}\varphi )S_t+rt\mathrm{sin}\theta \mathrm{cos}\varphi }{[tr\mathrm{sin}\theta \mathrm{cos}\varphi +S_t]S_t}}\right]_{R_1}^{R_2},`$ (84)
$`S_t=\sqrt{t^2+r^22rt\mathrm{cos}\varphi \mathrm{sin}\theta }.`$ (85)
The $`\varphi `$-integral is completely analytic, and yields
$`_\varphi (r,R_1,R_2,\theta )`$ $`=`$ $`\left[{\displaystyle \frac{4t}{rS_{}}}𝐊\left(\sqrt{{\displaystyle \frac{4tr\mathrm{sin}\theta }{S_{}^2}}}\right)\right]_{R_1}^{R_2},`$ (86)
$`=`$ $`\left[{\displaystyle \frac{4t}{rS_+}}𝐊\left(\sqrt{{\displaystyle \frac{4tr\mathrm{sin}\theta }{S_+^2}}}\right)\right]_{R_1}^{R_2},`$ (87)
$`S_\pm `$ $``$ $`\sqrt{t^2+r^2\pm 2rt\mathrm{sin}\theta }.`$ (88)
We thus have
$`a_{1/r^2}(r)`$ $`=`$ $`C_2g_{1/r^2}(r)`$ (89)
$`C_2`$ $`=`$ $`G\rho _2L^2/D=C_1/(8\delta )=2.5031C_1`$ (90)
$`g_{1/r^2}(r)`$ $`=`$ $`D{\displaystyle _{\pi /2\theta _0}^{\pi /2+\theta _0}}𝑑\theta \mathrm{sin}\theta _\varphi (r,R_1,R_2,\theta ).`$ (91)
This integral can be done numerically and is used in $`a_{1/r^2}(r)`$, shown in Figure 6. (The only numerical singularity problems are if both $`\theta =\pi /2`$ and also $`r`$ is either $`R_2`$ or $`R_1`$.)
The most interesting observation is that this result is very similar to that from the $`1/p`$-density cylindrical ring shown in Figure 4. (This point will be shown even better in Section VIII.) The fact that the density is falling off faster with distance ($`1/r^2`$ vs. $`1/p`$) is compensated for by the increasing spherical width, which is growing as $`r\mathrm{sin}\theta _0`$.
A comparison of the normalized acceleration, $`a_{1/r^2}(r)`$, with that for other models will be given in Section VIII.
### VII.2 Thin-ring configuration
To demonstrate the correctness of the assertion that the growing width of the wedge with distance caused the wedge to behave more like a $`1/p`$ ring, we now quickly look at the “thin-ring” $`1/p^2`$ problem. Keeping the same mass as before and using the formalism of Section III.2.2 yields (also see Eq. (44))
$`\rho _{T/p^2}(p)`$ $`=`$ $`{\displaystyle \frac{\rho _{T2}L^2}{p^2}},`$ (92)
$`\rho _{T2}`$ $`=`$ $`{\displaystyle \frac{(R_2R_1)}{L\mathrm{ln}(R_2/R_1)}}=(2.485)\rho _1`$ (93)
$`a_{T/p^2}(r)`$ $`=`$ $`C_{T2}g_{T/p^2}(r),`$ (94)
$`C_{T2}`$ $`=`$ $`C_1{\displaystyle \frac{(R_2R_1)}{4D\mathrm{ln}(R_2/R_1)}}=(24.85)C_1,`$ (95)
$`g_{T/p^2}(r)`$ $`=`$ $`D^2{\displaystyle _{R_1}^{R_2}}dp[{\displaystyle \frac{𝐊\left(\sqrt{\frac{4pr}{r^2+2pr+p^2}}\right)}{pr(r+p)}}`$ (96)
$`+{\displaystyle \frac{𝐄\left(\sqrt{\frac{4pr}{r^2+2pr+p^2}}\right)}{pr(rp)}}].`$
In Figure 7 we show a plot of $`a_{T/p^2}(r)`$. One clearly sees the difference between the $`1/r^2`$ wedge and the $`1/p^2`$ thin ring. With its rise going inward within the ring, $`a_{T/p^2}(r)`$ displays the higher mass concentration at $`r=R_1`$. (Again there is the thin-ring caveat that the spikes at $`r=\{R_1,R_2\}`$ would be finite cusps in a 3-d calculation.)
A comparison of the normalized acceleration, $`a_{T/p^2}(r)`$, with that for other models is also given in Section VIII.
## VIII Discussion
The different physical models we have investigated in this paper provide an intuitive understanding about what type of accelerations can be obtained from Kuiper Belt models. In particular, they can not easily yield a constant (or even an approximately constant) gravitational acceleration in a cylindrical system.
As to the specific gravitational accelerations in the plane of the ecliptic, $`a(r)`$, we found:
* Starting out with Figure 1, one sees that even a spherical shell with a $`1/r`$ density only yields an approximately constant acceleration near the outer edge of the shell.
* Continuing on to a “thin ring” with sharp edges, the $`1/p`$ density produces an acceleration that is singular at the edges of the ring and is approximately constant near the middle of the ring. (See Figure 2.)
* Contrary to this, the smoother-density, “thin-ring” Boss-Peale model produces a smooth acceleration at the inner edge and shows only a slight, narrow spike if the density has a small discontinuous jump at the outer edge instead of decreasing smoothly to infinity. (See Figure 3.) Thus, the physical differences in shape between the $`1/p`$ and Boss-Peale models end up being instructive.
* The full 3-dimensional, $`1/p`$ model, yields a finite acceleration everywhere, so the cusps at the edges of the ring are finite compared to the spikes of the 2-dimensional “thin-ring” approximation. (See Figure 4.)
* The 3-dimensional constant-density ring produces softer cusps yet a more undulatory variation than the $`1/p`$ ring. (See Figure 5.) It is intermediate in its effects between the 3-d, $`1/r`$ ring and the 2-d, Boss-Peale “thin” ring.
* The 3-dimensional wedge, with a spherical fall off in density of $`1/r^2`$, produces an acceleration that is very similar in shape to that from the $`1/p`$ cylindrical ring. (See Figure 6.) This is because the growing width of the wedge with distance approximately makes up for the faster fall off of density with distance.
* The above assertion is demonstrated by the contrasting behaviour of the $`1/p^2`$-density, thin ring’s $`a_{T/p^2}(r)`$, compared to the wedge’s $`a_{1/r^2}(r)`$. It varies much more in the belt and reaches a high maximum near $`r=R_1`$. (See Figure 7).
(We also mentioned how to extend these results to out of the plane of the ecliptic by taking $`r_z0`$ and then studying both $`a(r_x)`$ and $`a(r_z)=(d/dr_z)𝒱(𝐫)`$.)
The results emphasize how difficult it is to achieve a truly constant acceleration within a finite cylindrically-symmetric system (not even considering how much mass would be needed to mimic the Pioneer anomaly). This difficulty can be put in mathematical context. Consider just the “thin ring,” which is mathematically simpler than the full 3-d ring. Starting with Eq. (25), a constant acceleration between $`R_1`$ and $`R_2`$ would be produced by a density $`\rho _C(p)`$ that satisfied
$$r=\mathrm{Const}._{R_1}^{R_2}𝑑pp\rho _C(p)I_\varphi (r),$$
(97)
where $`I_\varphi (r)`$ is given in Eqs. (33) and (34). That is a complicated inverse problem. Formally it could be done by a decomposition into cylindrical harmonics, but that is not the point here.
Finally, in Figure 8, we show a direct comparison of the physical accelerations of (i) the 3-d, $`1/p`$ ring, (ii) the 2-d, Boss-Peale “thin ring,” (iii) the 3-d, constant ring, (iv) the 3-d, $`1/r^2`$ wedge, and (v) the ”thin,” $`1/p^2`$ ring, all with the same total mass, 1.96 $`M_{}`$. (As before we cut off the infinite spikes at the boundaries of the thin rings.) When $`r\mathrm{}`$, all the curves tend to $`[G_{ring}/r^2]`$, as they should. This is even though the differing density distributions produce quite different accelerations within the ring.
To within normalizations, the results in Figure 8 agree with the type of results published previously for Kuiper-Belt disks boss ; pioprd . Most importantly, within the ring the acceleration is not constant. Further, especially in the central portions of the rings, the accelerations are approximately two orders of magnitude too small to explain the Pioneer anomaly.
###### Acknowledgements.
I thank Claus Lämmerzahl and J. D. Anderson for their valuable comments, criticisms, and suggestions. The support of the U.S. DOE is acknowledged.
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# The Miura Map on the Line
## 1. Introduction and Main Results
The Miura map is the nonlinear mapping
(1.1)
$$B(r)=r^{}+r^2$$
which takes classical solutions of the modified Korteweg-de Vries (mKdV) equation to classical solutions of the Korteweg-de Vries (KdV) equation. More precisely, let
$$\mathrm{mKdV}(v):=v_t6v^2v_x+v_{xxx}$$
and
$$KdV(u):=u_t6uu_x+u_{xxx}$$
for functions $`v,u𝒞^{\mathrm{}}(\times )`$. Then
(1.2)
$$KdV(B(v))=\left(\mathrm{mKdV}(v)\right)_x+2v\mathrm{mKdV}(v)$$
so that $`KdV(B(v))=0`$ whenever $`mKdV(v)=0`$. The Miura map has been used extensively to relate existence and uniqueness results for solutions for the mKdV and KdV equations. More fundamentally, the global geometry of nonlinear differentiable mappings such as the Miura map has been studied to solve various nonlinear differential equations.
It is well-known that the KdV equation can be successfully studied with the help of the spectral theory of the Schrödinger operators
(1.3)
$$L_q:=d^2/dx^2+q.$$
In particular, in appropriate classes of potentials $`q`$, the spectrum of $`L_q`$ is preserved by the KdV flow applied to $`q`$. (See e.g. .) The Miura map is formally related with the Schrödinger operator as follows: the relation $`q=B(r)=r^{}+r^2`$ is equivalent to the following factorization of $`L_q`$:
(1.4)
$$L_q=L_{B(r)}=(_xr)^+(_xr)=(_xr)(_xr),$$
where $`_x=d/dx`$, and $`A^+`$ means the operator formally adjoint to an operator $`A`$ in functions on $``$ (with respect to the scalar product in $`L^2()`$).
The aim of this paper is to study the Miura map and its geometry on the real line with an emphasis on function spaces of low regularity. Our first result characterizes the range of the Miura map. We denote by $`B`$ the map (1.1) from real-valued functions in $`L_{\mathrm{loc}}^2()`$ into the local Sobolev space $`H_{\mathrm{loc}}^1()`$. If $`q`$ is a real-valued distribution in $`H_{\mathrm{loc}}^1()`$, the operator $`L_q`$ maps $`𝒞_0^{\mathrm{}}()`$ into the Sobolev space $`H^1()`$, so that $`(L_q\phi ,\phi )`$ is well-defined for any $`\phi 𝒞_0^{\mathrm{}}()`$. We write $`L_q0`$ if $`(L_q\phi ,\phi )0`$ for all $`\phi 𝒞_0^{\mathrm{}}()`$.
###### Theorem 1.1.
Let $`q`$ be a real-valued distribution belonging to $`H_{\mathrm{loc}}^1()`$. Then the following statements are equivalent.
(i) $`q\mathrm{Im}(B)`$, i.e., $`q=r^{}+r^2`$ for some real-valued function $`rL_{\mathrm{loc}}^2()`$.
(ii) The equation $`L_qy=0`$ has a strictly positive solution $`yH_{\mathrm{loc}}^1()`$.
(iii) $`L_q0`$.
Next, we consider the restriction of $`B`$ to the Sobolev space $`H^\beta ()`$ for $`\beta 0`$; we denote this restriction by $`B_\beta `$. Although $`B_\beta `$ has range contained in $`H^{\beta 1}()`$, it is not true that $`\mathrm{Im}(B_\beta )=\mathrm{Im}(B)H^{\beta 1}()`$; rather, an additional condition is needed to characterize its range.
###### Theorem 1.2.
Let $`\beta 0`$ be arbitrary. A real-valued distribution $`qH^{\beta 1}()`$ belongs to $`\mathrm{Im}(B_\beta )`$ if and only if
(i) $`L_q0`$, and
(ii) $`q`$ can be presented in the form $`q=f^{}+g`$ for $`fL^2()`$ and $`gL^1\left(\right)`$.
In addition, we will give an alternative characterization of $`\mathrm{Im}(B_\beta )`$ in terms of a “special integral” of $`q`$ on $``$ which coincides with the ordinary integral of $`q`$ if $`qL^1\left(\right)`$ (see Theorem 4.5).
We also study the geometry of the Miura map. Whereas the Miura map on spaces of periodic functions is known to be a global fold, the situation for non-periodic functions is completely different. Let
(1.5)
$$E_1=\{q\mathrm{Im}(B):B^1(q)\text{ consists of a single point}\}$$
and
(1.6)
$$E_2=\{q\mathrm{Im}(B):B^1(q)\text{ is homeomorphic to an interval}\}.$$
Here we consider $`B^1(q)`$ with the natural Fréchet topology of $`L_{\mathrm{loc}}^2()`$.
###### Theorem 1.3.
$`\mathrm{Im}(B)=E_1E_2`$ and both $`E_1`$ and $`E_2`$ are dense in $`\mathrm{Im}(B)`$ in the natural Fréchet topology of $`H_{\mathrm{loc}}^1()`$.
As an application of our results on the Miura map, we prove an existence result for solutions of the Korteweg-de Vries equation in $`H^1()`$ for initial data in the range $`\mathrm{Im}(B_0)`$ of the Miura map $`B_0:L^2()H^1()`$. We follow the approach of Tsutsumi , who proved such an existence result when the initial data is a positive, finite Radon measure on $``$. His arguments, combined with our results on the Miura map $`B_0,`$ lead to the following theorem.
###### Theorem 1.4.
Assume that $`u_0\mathrm{Im}(B_0)`$. Then there exists a global weak solution of KdV with $`u(t)\mathrm{Im}\left(B_0\right)`$ for all $`t`$. More precisely:
(i) $`uL^{\mathrm{}}(,H^1())L_{\mathrm{loc}}^2(^2)`$;
(ii) for all functions $`\phi 𝒞_0^{\mathrm{}}(^2)`$, the identity
$$_{}_{}\left(u\phi _tu\phi _{xxx}+3u^2\phi _x\right)𝑑x𝑑t=0$$
holds, and
(iii) $`lim_{t0}u(t)=u_0`$ in $`H^1()`$.
We state and prove a slightly stronger version of the above theorem in Section 6.
In Appendix C, we provide a few comments on the Miura map as well as on other work related to the results presented in this paper.
Acknowledgements. Peter Perry and Mikhail Shubin thank the Mathematical Institute at the University of Zürich for its hospitality during part of the time that this work was done. The authors thank Carlos Tomei for bringing the paper of McKean and Scovel to their attention. They also thank Vladimir Maz’ya and Tatyana Shaposhnikova for pointing out the reference (see the proof of Lemma 2.2).
## 2. Preliminaries
### 2.1. Spaces of Distributions
For $`\alpha `$, we denote by $`H^\alpha ()`$ the completion of $`𝒞_0^{\mathrm{}}()`$ in the norm
$$\phi _\alpha =\left(\left(1+\left|\xi \right|^2\right)^\alpha \left|\widehat{u}(\xi )\right|^2\frac{d\xi }{2\pi }\right)^{1/2}$$
where
$$\widehat{u}(\xi )=\mathrm{exp}\left(i\xi x\right)u(x)𝑑x.$$
Clearly, $`H^0()=L^2()`$ and $`H^\alpha ()H^\beta ()`$ if $`\alpha \beta `$. Therefore $`H^\alpha ()L^2()`$ if $`\alpha 0`$. For any $`\alpha `$, $`H^\alpha ()`$ is a space of tempered distributions. A distribution $`u`$ belongs to $`H_{\mathrm{loc}}^\alpha ()`$ if $`\chi uH^\alpha ()`$ for any function $`\chi 𝒞_0^{\mathrm{}}()`$. We consider $`H_{\mathrm{loc}}^\alpha ()`$ and $`L_{\mathrm{loc}}^p()`$ with $`p1`$ with their natural Fréchet topology - see , chapter 31-12.
A classical result of distribution theory (see, for example, chapter 1 of ) asserts that if $`qH_{\mathrm{loc}}^1()`$, then $`q=Q^{}`$ for a function $`QL_{\mathrm{loc}}^2()`$. If $`qH^{\beta 1}()`$ for $`\beta 0`$ we have a sharper result. Let
$$H^{\mathrm{}}()=\underset{\beta 0}{}H^\beta ()𝒞^{\mathrm{}}()$$
###### Lemma 2.1.
Let $`\beta 0`$ and let $`qH^{\beta 1}()`$. There exist functions $`fH^\beta ()`$ and $`gH^{\mathrm{}}()`$ so that $`q=f^{}+g`$ as elements of $`H^{\beta 1}()`$.
###### Proof.
Let $`\psi 𝒞_0^{\mathrm{}}()`$ with $`\psi (\xi )=1`$ near $`\xi =0`$, and choose
$$\widehat{f}(\xi )=\left(i\xi \right)^1\left(1\psi (\xi )\right)\widehat{q}(\xi )$$
and
$$\widehat{g}(\xi )=\psi (\xi )\widehat{q}(\xi ).$$
Finally, the following result will be useful in studying the continuity of the Miura map and the regularity of solutions to the Riccati equation $`q=r^{}+r^2`$.
###### Lemma 2.2.
The multiplication $`\{u,v\}uv`$ can be extended from the bilinear map $`𝒞_0^{\mathrm{}}()\times 𝒞_0^{\mathrm{}}()𝒞_0^{\mathrm{}}()`$ to continuous bilinear maps
(2.1)
$$H^\beta ()\times H^\beta ()H^\beta (),\beta >1/2,$$
(2.2)
$$H^{1/2}()\times H^{1/2}()H^{1/2\delta }()\text{ for any }\delta >0,$$
(2.3)
$$H^\beta ()\times H^\beta ()H^{2\beta 1/2}(),0<\beta <1/2,$$
(2.4)
$$L^2()\times L^2()L^1()H^{1/2\delta }()\text{ for any }\delta >0\text{.}$$
###### Proof.
All of the statements, except the last one, are particular cases of more general multidimensional results formulated, for example, in Theorem 1 of section 4.6.1 of . The last statement is obvious except the last inclusion; this follows by duality from the imbedding $`H^{1/2+\delta }()L^{\mathrm{}}()`$, which is a particular case of a well-known Sobolev imbedding theorem. ∎
Let us also introduce a notation $`(,)`$ for miscellaneous sesquilinear pairings extending the pairing
$$(u,v)=_{}u(x)\overline{v(x)}𝑑x,u,v𝒞_0^{\mathrm{}}(),$$
by continuity. In particular, we will use the extended pairings in the following cases:
(i) $`u`$ is a distribution on $``$, $`v𝒞_0^{\mathrm{}}()`$;
(ii) $`uH^s()`$, $`vH^s()`$, where $`s`$;
(iii) $`uH_{\mathrm{l}oc}^s()`$, $`vH_{\mathrm{c}omp}^s()`$, where $`s`$ and $`H_{comp}^s()`$ is the space of compactly supported distributions from $`H^s()`$.
The integration by parts formula
(2.5)
$$(U^{},v)=(U,v^{})$$
holds in all these cases, e.g. for $`UH_{\mathrm{l}oc}^{s+1}()`$, $`vH_{\mathrm{c}omp}^s()`$, $`s`$, by continuity of the pairing.
### 2.2. Continuity of the Miura Map
It is easy to see that the Miura map defines a bounded continuous map from $`L^2()`$ into $`H^1()`$ (we use the standard embedding $`L^1()H^1()`$ which follows by duality from the Sobolev embedding $`H^1()L^{\mathrm{}}()`$). Localizing this, we see that the Miura map defines a bounded continuous mapping from $`L_{\mathrm{loc}}^2()`$ into $`H_{\mathrm{loc}}^1()`$. This is extended to more general Sobolev spaces as follows.
###### Proposition 2.3.
Let $`\beta 0`$. The Miura map is a continuous mapping from $`H^\beta ()`$ into $`H^{\beta 1}()`$ and from $`H_{\mathrm{loc}}^\beta ()`$ into $`H_{\mathrm{loc}}^{\beta 1}()`$ which is bounded, i.e. maps bounded subsets into bounded subsets.
###### Remark 2.4.
For the definition of bounded sets in Fréchet spaces see e.g. , Chapter 14, especially Proposition 14.5.
###### Proof.
It suffices to prove the first statement since the second follows by localization. It is clear that the map $`rr^{}`$ is a bounded continuous map from $`H^\beta ()`$ into $`H^{\beta 1}()`$ so it suffices to show that the map $`rr^2`$ is a bounded continuous map from $`H^\beta ()`$ into $`H^{\beta 1}()`$. This follows from Lemma 2.2 and the trivial embedding of $`H^\alpha ()`$ into $`H^\gamma ()`$ if $`\alpha \gamma `$. ∎
### 2.3. A First-Order System
Let $`yH_{\mathrm{loc}}^1()`$ be a solution of the equation
(2.6)
$$y^{\prime \prime }+qy=0$$
where $`qH_{\mathrm{loc}}^1()`$. Then we can conveniently rewrite the equation in the form of a first-order system
(2.7)
$$\{\begin{array}{ccc}y^{}& =& Qy+u\hfill \\ u^{}& =& Q^2yQu\hfill \end{array}$$
where $`QL_{\mathrm{loc}}^2()`$ is a real-valued function with $`Q^{}=q`$ (this is a well-known procedure in the study of differential operators with singular coefficients; see §1.1 of and references therein). This system is equivalent to equation (2.6) in the following sense. If $`yH_{\mathrm{loc}}^1()`$ is a solution of (2.6), then taking $`u=y^{}Qy`$ we obtain by straightforward substitution that the pair $`\{y,u\}`$ satisfies (2.7) in the sense of distributions. It follows that $`uW_{\mathrm{loc}}^{1,1}()`$, the space of $`L_{\mathrm{loc}}^1()`$-functions with distributional derivatives in $`L_{\mathrm{loc}}^1()`$. On the other hand, if $`y`$ and $`u`$ belong to $`W_{\mathrm{loc}}^{1,1}()`$ and the pair $`\{y,u\}`$ satisfies (2.7), then in fact $`yH_{\mathrm{loc}}^1()`$ and $`y`$ satisfies (2.6).
Since the coefficients of the linear system (2.7) are in $`L_{\mathrm{loc}}^1()`$, the standard existence and uniqueness result holds for the corresponding initial value problem on the whole real line, and the solutions $`\{y,u\}`$ depend continuously on the initial data.
The next lemma shows that nonnegative solutions of $`L_qy=0`$ are either strictly positive or identically zero.
###### Lemma 2.5.
Suppose that $`yH_{\mathrm{loc}}^1(I)`$ is a solution of $`L_qy=0`$ with a real-valued $`qH_{\mathrm{loc}}^1(I)`$, where $`I`$ is an open interval in $``$. Assume that $`y(x_0)=0`$ for some $`x_0I`$. Denote $`u=y^{}Qy`$ as in (2.7). Then the following statements hold true:
A. We have the following trichotomy for the behavior of $`y`$ near $`x_0`$:
(i) If $`u(x_0)=0`$, then $`y0`$ on $`I`$.
(ii) If $`u(x_0)>0`$ then in a neighborhood of $`x_0`$ we have $`y(x)<0`$ for $`x<x_0`$, and $`y(x_0)>0`$ for $`x>x_0`$.
(iii) If $`u(x_0)<0`$ then in a neighborhood of $`x_0`$ we have $`y(x)>0`$ for $`x<x_0`$, and $`y(x_0)<0`$ for $`x>x_0`$.
B. If $`y0`$, then all zeros of $`y`$ on $`I`$ are isolated.
C. If $`y(x)0`$ near $`x_0`$ then $`y(x)0`$ on $`I`$.
###### Proof.
Clearly, $`(i)`$ follows from the uniqueness of solution $`\{y,u\}`$ of (2.7) with the given initial conditions $`y(x_0)=y_0,u(x_0)=u_0`$.
Define $`z(x)=y(x)\mathrm{exp}\left(_{x_0}^xQ(s)𝑑s\right)`$ and note that, as $`u`$ is absolutely continuous, so is $`z^{}(x)=u(x)\mathrm{exp}\left(_{x_0}^xQ(s)𝑑s\right)`$. Clearly, $`z(x_0)=0`$, and $`z(x)`$ has the same sign as $`y(x)`$ for all $`xI`$. Also, $`u(x_0)>0`$ is equivalent to $`z^{}(x_0)>0`$, so $`(ii)`$ and $`(iii)`$ immediately follow.
Clearly, $`B`$ and $`C`$ follow from $`A`$. ∎
## 3. Positivity, Positive Solutions, and the Miura Map
In this section we prove Theorem 1.1. First, we describe the connection between the Miura map and positive solutions of $`L_qy=0`$. For a real-valued distribution $`qH_{\mathrm{loc}}^1()`$, let $`Pos(q)`$ denote the (possibly empty) set of functions $`yH_{\mathrm{loc}}^1()`$ with the properties that $`L_qy=0`$, $`y(x)>0`$ for all $`x`$, and $`y(0)=1`$.
###### Lemma 3.1.
Let $`q`$ be a real-valued distribution belonging to $`H_{\mathrm{loc}}^1()`$.
(a) If $`yPos(q)`$ then $`q=B(y^{}/y)`$.
(b) If $`q=B(r)`$ for some $`rL_{\mathrm{loc}}^2()`$ then $`y(x)=\mathrm{exp}\left(_0^xr(s)𝑑s\right)`$ belongs to $`Pos(q)`$.
The proof is easily obtained by straightforward calculations.
The maps
$$Pos(q)y\frac{d}{dx}\mathrm{log}(y(x))B^1(q)$$
and
$$B^1(q)r\mathrm{exp}\left(_0^xr(s)𝑑s\right)Pos(q)$$
are continuous if we topologize $`B^1(q)`$ with the topology induced from $`L_{\mathrm{loc}}^2()`$ and $`Pos(q)`$ with the topology induced from $`H_{\mathrm{loc}}^1()`$. These maps are mutual inverses. Hence, we have shown:
###### Proposition 3.2.
The set $`B^1(q)`$ is nonempty if and only if $`Pos(q)`$ is nonempty. For any $`rB^1(q)`$, $`B^1(B(r))`$ is homeomorphic to $`Pos(B(r))`$.
Next, we show that if $`L_q0`$, then $`Pos(q)`$ is nonempty. To this end, we introduce the sesquilinear forms
(3.1)
$$𝔱_q(\phi ,\psi )=_{}\phi ^{}(x)\overline{\psi ^{}}(x)𝑑x+(q,\overline{\phi }\psi )$$
and
(3.2)
$$𝔱_{q,I}(\phi ,\psi )=_I\phi ^{}(x)\overline{\psi ^{}}(x)𝑑x+(q,\overline{\phi }\psi )$$
defined respectively on $`𝒞_0^{\mathrm{}}()`$ and $`𝒞_0^{\mathrm{}}(I)`$, where $`I=(a,b)`$ is a bounded, open interval of $``$.
The form $`𝔱_q`$ is also well defined by (3.1) for $`\phi ,\psi H_{comp}^1()`$, where $`H_{comp}^1()`$ is the space of compactly supported functions from $`H^1()`$. Note that $`𝔱_q(\phi ,\phi )=(L_q\phi ,\phi )`$ if $`\phi 𝒞_0^{\mathrm{}}()`$, so that if $`L_q0`$, then both $`𝔱_q`$ and $`𝔱_{q,I}`$ are positive quadratic forms. Approximating $`\phi H_{comp}^1()`$ by functions from $`𝒞_0^{\mathrm{}}()`$, we easily obtain that $`𝔱_q(\phi ,\phi )0`$ for all $`\phi H_{comp}^1()`$ as well.
It is easy to see that $`𝔱_{q,I}`$ admits a closure, which has the domain
$$H_0^1(I)=\{\psi H^1(I):\psi (a)=\psi (b)=0\}.$$
(See also Lemma 1.8 of .) It is a closed positive quadratic form which will also be denoted $`𝔱_{q,I}`$. Note that if $`\phi ,\psi H_0^1(I)`$ and $`\phi _0,\psi _0`$ are their extensions on $``$ by $`0`$, then $`\phi _0,\psi _0H_{comp}^1()`$, and
(3.3)
$$𝔱_{q,I}(\phi ,\psi )=𝔱_q(\phi _0,\psi _0).$$
Let $`L_{q,I}`$ be the self-adjoint operator associated to $`𝔱_{q,I}`$ by the Friedrichs construction. Clearly, $`L_{q,I}`$ has positive spectrum. Moreover, it has compact resolvent, or, equivalently, discrete spectrum. (This follows from the compactness of the imbedding of $`H_0^1(I)`$ into $`L^2(I)`$; see also , where the asymptotics of the eigenvalues is found.) By the min-max principle, the lowest eigenvalue of $`L_{q,I}`$ is given by
$$\lambda _0(I)=inf\{𝔱_{q,I}(\phi ,\phi ):\phi H_0^1(I)\text{ and }\phi _{L^2(I)}=1\}0$$
and the infimum is achieved by a corresponding eigenfunction $`hH_0^1(I)`$.
###### Lemma 3.3.
Let $`q`$ be a real-valued distribution belonging to $`H_{\mathrm{loc}}^1()`$. If $`L_q0`$, then $`\lambda _0(I)>0`$ for every bounded open interval $`I`$.
###### Proof.
As $`\lambda _0(I)0`$ it suffices to show that no bounded interval $`I`$ has $`\lambda _0(I)=0`$. Suppose, on the contrary, that such an interval $`I`$ exists, and let $`hH_0^1(I)`$ be an $`L^2`$-normalized, real-valued eigenfunction with the eigenvalue $`0`$, so, in particular, $`𝔱_{q,I}(h,h)=0`$. Extend $`h`$ to a function $`\eta =h_0H_{\mathrm{c}omp}^1()`$ as above, i.e. by setting $`\eta (x)=0`$ if $`xI`$. By (3.3), we conclude that $`𝔱_q(\eta ,\eta )=𝔱_{q,I}(h,h),`$ hence $`𝔱_q(\eta ,\eta )=0`$, So, for any $`\phi 𝒞_0^{\mathrm{}}()`$ and $`t`$,
$$𝔱_q(\eta +t\phi ,\eta +t\phi )=2t\mathrm{Re}𝔱_q(\eta ,\phi )+t^2\text{ }𝔱_q(\phi ,\phi )$$
is nonnegative due to positivity of $`𝔱_q`$. It follows that $`\mathrm{Re}𝔱_q(\eta ,\phi )=0`$ for all $`\phi 𝒞_0^{\mathrm{}}()`$, hence $`𝔱_q(\eta ,\phi )=0`$ also for all $`\phi 𝒞_0^{\mathrm{}}()`$. It follows that $`\eta `$ solves the equation $`L_q\eta =0`$ and has a compact suppport. By Lemma 2.5 we obtain $`\eta (x)=0`$ identically. This contradicts the assumption that $`h_{L^2(I)}=1`$, and the lemma is proved. ∎
###### Corollary 3.4.
Let $`q`$ be a real-valued distribution from $`H_{\mathrm{loc}}^1()`$. If $`L_q0`$ and $`yH_{\mathrm{loc}}^1()`$ solves $`L_qy=0`$, then $`y`$ can have at most one zero.
###### Proof.
If $`y(a)=y(b)=0`$ for $`a<b`$, then zero is a Dirichlet eigenvalue of $`L_{q,I}`$ with $`I=(a,b)`$, contradicting Lemma 3.3.
###### Proposition 3.5.
Let $`q`$ be a real-valued distribution which belongs to $`H_{\mathrm{loc}}^1()`$. If $`L_q0`$, then $`L_qy=0`$ has a strictly positive solution.
###### Proof.
For $`c\backslash \left\{0\right\}`$, let $`\{y,u\}H_{\mathrm{loc}}^1()\times W_{\mathrm{loc}}^{1,1}()`$ be the unique solution of (2.7) with $`y(c)=0`$ and $`u(c)=1`$. As $`y`$ satisfies $`L_qy=0`$, Corollary 3.4 implies that $`y(0)0`$. Denote by $`\{y_c,u_c\}`$ the scaled solution of (2.7),
(3.4)
$$y_c(x)=y(x)/y(0),u_c(x)=u(x)/y(0).$$
Then $`y_cH_{\mathrm{loc}}^1()`$ is the unique solution of $`L_qy=0`$ with $`y(0)=1`$ and $`y(c)=0`$. Next, choose $`c,c^{}`$ so that $`0<c^{}<c`$. Note that $`w(x)=y_c(x)y_c^{}(x)`$ is a solution of $`L_qy=0`$ with $`w(0)=0`$, but $`w(c)>0`$. By Corollary 3.4, $`w(x)>0`$ for $`x>0`$ and by Lemma 2.5 and Corollary 3.4 $`w(x)<0`$ for $`x<0`$. It follows that, on the half-line $`c>0`$, the map $`cy_c(x)`$ is monotone decreasing for any given $`x<0`$ and monotone increasing for any given $`x>0`$.
We wish to construct a positive solution of $`L_qy=0`$ by taking the limit $`c+\mathrm{}`$. To this end denote by $`(\stackrel{~}{y}_1,\stackrel{~}{u}_1)`$ and $`(\stackrel{~}{y}_2,\stackrel{~}{u}_2)`$ the fundamental solutions of (2.7), determined by $`\stackrel{~}{y}_1(0)=1`$, $`\stackrel{~}{u}_1(0)=0`$ and $`\stackrel{~}{y}_2(0)=0,\stackrel{~}{u}_2(0)=1`$ respectively. By Lemma 3.3, $`\stackrel{~}{y}_2(1)0.`$ Hence, for any $`\alpha `$
$$z(x;\alpha ):=\stackrel{~}{y}_1(x)+\frac{\alpha \stackrel{~}{y}_1(1)}{\stackrel{~}{y}_2(1)}\stackrel{~}{y}_2(x)$$
is the unique solution of $`L_qy=0`$ with $`y(1)=\alpha `$ and $`y(0)=1.`$ Moreover, it follows that for any $`x,`$ $`\alpha z(x;\alpha )`$ is continuous. (In fact, the map $`\alpha z(;\alpha )`$ is an affine, hence continuous map $`H_{\mathrm{l}oc}^1()`$.)
Note that $`y_c(x)=z(x;y_c(1))`$. Now consider the solution $`y_n`$ of $`L_qy=0`$, $`y(0)=1`$, $`y(n)=0`$. Then $`\alpha _n=y_n(1)`$ is a strictly positive, decreasing sequence. Let $`\alpha _{\mathrm{}}=`$ $`lim_n\mathrm{}\alpha _n`$. By the continuity of $`z(x;\alpha )`$ with respect to $`\alpha `$ and the fact that $`y_n(x)=z(x;\alpha _n),`$ it then follows that $`z(x;\alpha _{\mathrm{}})=lim_n\mathrm{}y_n(x)`$ for any $`x`$. We claim that $`\alpha _{\mathrm{}}>0`$ and that $`z(x;\alpha _{\mathrm{}})`$ is positive. To prove this, first note that for any $`n1,`$ $`y_n(x)0`$ for all $`xn`$ and hence $`lim_n\mathrm{}y_n(x)0`$ for all $`x`$. If $`\alpha _{\mathrm{}}=0`$, then
$$0\underset{n\mathrm{}}{lim}y_n(2)=\underset{n\mathrm{}}{lim}z(2;\alpha _n)=z(2;0)=y_1(2)<0,$$
a contradiction. Hence $`\alpha _{\mathrm{}}>0`$. As a consequence, $`z(x;\alpha _{\mathrm{}})`$ is a nonnegative solution of $`L_qy=0`$ and hence strictly positive by Lemma 2.5. ∎
*Proof of Theorem 1.1*. In the statement of Theorem 1.1, we have (ii)$``$(i) by Lemma 3.1(a). To show that (i)$``$(iii), we compute that, for $`q=r^{}+r^2`$ and any $`\phi 𝒞_0^{\mathrm{}}()`$,
(3.5)
$$(L_q\phi ,\phi )=\left|\phi ^{}r\phi \right|^2𝑑x0.$$
Finally, (iii)$``$(ii) by Proposition 3.5. $`\mathrm{}`$
## 4. The Image of the Miura Map
In this section we prove Theorem 1.2. Recall that $`B_\beta `$ denotes the restriction of the Miura map to the Sobolev space $`H^\beta ()`$ for $`\beta 0`$.
We begin by considering the case $`\beta =0`$. In the light of Theorem 1.1, we need to find necessary and sufficient conditions on a potential $`qH^1()`$ so that there exists a solution $`rL^2()`$ of the Riccati equation $`r^{}+r^2=q`$. Hartman , chapter XI.7, Lemma 7.1 has studied this problem for continuous $`q`$ and his arguments still apply in our more general setting.
###### Lemma 4.1.
Suppose that $`qH^1()`$ is a real-valued distribution and that
(4.1)
$$\underset{\left|T\right|>1}{sup}\left|\frac{1}{T}_0^TQ(x)𝑑x\right|<+\mathrm{}$$
for an antiderivative $`QL_{\mathrm{loc}}^2()`$ of $`q`$. Then every solution $`rL_{\mathrm{loc}}^2()`$ of the Riccati equation $`r^{}+r^2=q`$ belongs to $`L^2()`$. Conversely, if $`rL^2()`$, then every antiderivative $`Q`$ of $`q=r^{}+r^2`$ satisfies (4.1).
###### Proof.
(i) Assume that $`QL_{\mathrm{loc}}^2()`$, $`Q^{}=q`$, and $`Q`$ satisfies (4.1). We need to show that for any solution $`rL_{\mathrm{loc}}^2()`$ of $`r^{}+r^2=q`$, the integrals $`_0^{\mathrm{}}r^2(s)𝑑s`$ and $`_{\mathrm{}}^0r^2(s)𝑑s`$ are finite. Let us show that the first integral is finite. Since $`Q`$ is an antiderivative of $`q`$, it follows that, for a constant $`C`$,
(4.2)
$$r(x)+_0^xr^2(s)𝑑s=Q(x)+C.$$
By assumption, the Césaro mean $`T^1_0^T()𝑑x`$ of the right-hand side of (4.2) is bounded as $`T+\mathrm{}`$. We suppose that $`_0^xr^2(s)𝑑s+\mathrm{}`$ as $`x+\mathrm{}`$ and obtain a contradiction as follows.
First note that if $`fL_{\mathrm{loc}}^1()`$ and $`f(x)+\mathrm{}`$ as $`x+\mathrm{}`$ then the same holds for its Césaro mean, i.e.
(4.3)
$$\frac{1}{T}_0^Tf(x)𝑑x+\mathrm{}\mathrm{as}T+\mathrm{}.$$
Hence, taking Césaro means of (4.2) we see that, if $`_0^xr^2(s)𝑑s+\mathrm{}`$ as $`x+\mathrm{}`$, then $`T^1_0^Tr(x)𝑑x\mathrm{}`$ as $`x+\mathrm{}`$. Moreover, there is a $`T_0`$ so that for all $`T>T_0`$,
$$\frac{2}{T}_0^Tr(x)𝑑x\frac{1}{T}_0^T\left(_0^xr^2(s)𝑑s\right)𝑑x>0.$$
By the Cauchy-Schwarz inequality,
$$T^1_0^Tr(x)𝑑xT^{1/2}\left(_0^Tr^2(x)𝑑x\right)^{1/2}$$
so that
$$4T_0^Tr^2(x)𝑑x\left[_0^T\left(_0^xr^2(s)𝑑s\right)𝑑x\right]^2.$$
Setting $`I(T)=_0^T_0^xr^2(s)𝑑s𝑑x`$, we have that
$$4TI^{}(T)I(T)^2$$
from which it follows by integration that
$$\frac{1}{I(T_0)}\frac{1}{I(T)}\frac{1}{4}\mathrm{log}\left(T/T_0\right).$$
This contradicts that, by (4.3), $`I(T)+\mathrm{}`$ as $`T+\mathrm{}`$. A similar argument shows that $`_{\mathrm{}}^0r^2(s)𝑑s`$ is finite. Hence $`rL^2()`$.
(ii) If, on the other hand, $`q=r^{}+r^2`$ for $`rL^2()`$, then the function $`Q(x)=r(x)+_0^xr^2(s)𝑑s`$ satisfies (4.1) by the Cauchy-Schwarz inequality applied in $`[0,T]`$. ∎
###### Corollary 4.2.
Suppose that $`q\mathrm{Im}(B_0)`$, i.e., $`q=r^{}+r^2`$ for some $`rL^2()`$. If $`uL_{\mathrm{loc}}^2()`$ solves the Riccati equation $`u^{}+u^2=q`$, then $`uL^2()`$.
###### Proposition 4.3.
A real-valued distribution $`qH^1()`$ belongs to $`\mathrm{Im}(B_0)`$ if and only if
(i) $`L_q0`$, and
(ii) $`q`$ can be presented as $`q=f^{}+g`$ for real-valued functions $`fL^2()`$ and $`gL^1()`$.
###### Proof.
Suppose that $`q\mathrm{Im}(B_0)`$, i.e., $`q=r^{}+r^2`$ for some $`rL^2()`$. Then $`L_q0`$ by Theorem 1.1 and $`q=f^{}+g`$ with $`f=r`$, $`g=r^2`$. On the other hand, suppose that $`qH^1()`$ with $`L_q0`$, and $`q=f^{}+g`$ for $`fL^2()`$ and $`gL^1()`$. By Theorem 1.1, $`q\mathrm{Im}(B)`$, so $`q=r^{}+r^2`$ for some $`rL_{\mathrm{loc}}^2()`$. The antiderivative $`Q(x)=f(x)+_0^xg(s)𝑑s`$ obeys the condition (4.1), so $`rL^2()`$ by Lemma 4.1. ∎
###### Proposition 4.4.
The set $`\mathrm{Im}(B_0)`$ has no interior points, and hence is not open in $`H^1()`$. Further, the set $`\mathrm{Im}(B_0)`$ is not closed in $`H^1()`$.
###### Proof.
First we show that $`\mathrm{Im}(B_0)`$ has empty interior. If $`q𝒞_0^{\mathrm{}}()\mathrm{Im}(B_0)`$, we can perturb $`q`$ by a small potential well far separated from the support of $`q`$ and create a bound state. More precisely, for $`\epsilon >0`$, let
$$v_\epsilon (x)=\{\begin{array}{ccc}\epsilon & & \left|x\right|<1/(2\epsilon )\hfill \\ 0& & \left|x\right|1/(2\epsilon )\hfill \end{array}.$$
Observe that $`v_\epsilon (x)𝑑x=1`$ but $`v_ϵ_{L^2()}=\epsilon `$. Suppose that $`q𝒞_0^{\mathrm{}}()\mathrm{Im}(B_0)`$ with support contained in $`[a,a]`$. Let $`w_\epsilon (x)=v_\epsilon (x2a2\epsilon ^1)`$; then $`w_\epsilon `$ has support disjoint from the one of $`q`$. By choosing $`\epsilon `$ sufficiently small, we can assure that the potential $`q_\epsilon =q+w_\epsilon `$ is close to $`q`$ in $`L^2()`$ norm. Let $`\chi 𝒞_0^{\mathrm{}}()`$ be a nonnegative function with $`\chi (x)=1`$ for $`\left|x\right|<1/(2\epsilon )`$, $`\chi (x)=0`$ for $`\left|x\right|>\epsilon ^1`$, and $`\left|\chi ^{}(x)\right|3\epsilon `$. Finally, let $`\eta (x)=`$ $`\chi \left(x2a2\epsilon ^1\right)`$. Then
$`(L_{q_\epsilon }\eta ,\eta )`$ $`={\displaystyle \left|\eta ^{}(x)\right|^2}118\epsilon 1.`$
Hence, by Theorem 1.1, $`q_\epsilon \mathrm{Im}(B)`$ for $`0<\epsilon <1/18`$. Since $`𝒞_0^{\mathrm{}}()`$ is norm-dense in $`H^1()`$, this shows that $`\mathrm{Im}(B_0)`$ contains no open neighborhood in the norm topology of $`H^1()`$.
Next, we show that $`\mathrm{Im}(B_0)`$ is not closed. Suppose that $`q`$ is any nonnegative function with $`qL^2()`$ but $`qL^1()`$. We may approximate $`q`$ by nonnegative potentials $`q_k𝒞_0^{\mathrm{}}()`$ so that $`q_kq`$ in $`L^2()`$. Hence, $`q_kq`$ in $`H^1()`$ and by Proposition 4.3, $`q_k\mathrm{Im}(B_0)`$ for any $`k1.`$ Moreover, since $`L_q0`$, it follows from Theorem 1.1 that $`q\mathrm{Im}(B)`$. On the other hand, as $`q0,`$ no antiderivative $`Q`$ of $`q`$ satisfies condition (4.1). Thus $`\mathrm{Im}(B_0)`$ is not closed in the norm topology on $`H^1()`$. ∎
The image of $`B_0`$ can also be characterized by a ‘special integral’ of $`q`$. Let $`\left\{\chi _n\right\}_{n1}`$ be a sequence of nonnegative $`𝒞_0^{\mathrm{}}()`$ functions with (i) $`\chi _n(x)=1`$ for $`\left|x\right|n`$, (ii) $`\chi _n(x)=0`$ for $`\left|x\right|n+1`$, and (iii) $`\left|\chi _n^{}(x)\right|2`$ for all $`x`$. Given $`qH^1()`$, we define the special integral of $`q`$, denoted $`\left[q\right]`$, to be the number $`lim_n\mathrm{}(q,\chi _n)`$ if this limit exists and is finite. One easily checks that $`\left[q\right]`$ is well-defined, i.e., does not depend on the choice of sequence $`\left\{\chi _n\right\}_{n1}`$ satisfying properties (i), (ii), and (iii) above. If $`q\mathrm{Im}(B_0)`$ then, for any $`rB_0^1(q)`$,
(4.4) $`\underset{n\mathrm{}}{lim}(q,\chi _n)`$ $`=\underset{n\mathrm{}}{lim}\left\{(r,\chi _n^{})+(r^2,\chi _n)\right\}=r_{L^2()}^2`$
which shows that $`\left[q\right]0`$ on $`\mathrm{Im}(B_0)`$ with $`\left[q\right]=0`$ if and only if $`q=0`$. Moreover, if $`r_1`$ and $`r_2`$ belong to $`B^1(q)`$, then $`r_1_{L^2()}=r_2_{L^2()}`$.
The special integral has the following properties:
1. $`\left[f\right]=_{}f𝑑x`$ if $`fL^1()H^1()`$;
2. $`Dom\left([]\right)`$ is a linear subspace in $`H^1()`$, and $`f\left[f\right]`$ is linear;
3. $`\left[f^{}\right]=0`$ for any $`fL^2()`$;
4. If $`fL^2()`$, then $`\left[f\right]`$ exists if and only if $`f`$ is conditionally integrable, i.e. the limit $`lim_T\mathrm{}_T^Tf(x)𝑑x`$ exists. In this case, the limit equals $`\left[f\right].`$
Using the special integral we can give an alternative characterization of $`\mathrm{Im}(B_0)`$.
###### Theorem 4.5.
A real-valued distribution $`qH^1()`$ belongs to $`\mathrm{Im}(B_0)`$ if and only if:
(i) $`L_q0`$, and
(ii) $`\left[q\right]`$ exists.
Moreover, for any $`q\mathrm{Im}(B_0),`$ one has $`\left[q\right]0.`$
###### Proof.
First, suppose that $`L_q0`$ and $`\left[q\right]`$ exists. To prove that $`q\mathrm{Im}(B_0)`$, it suffices by Lemma 4.1 to show that $`q`$ has an antiderivative $`Q`$ with bounded Césaro means. By Lemma 2.1, any $`qH^1()`$ may be written $`q=f^{}+g`$ for $`f`$ and $`g`$ belonging to $`L^2()`$. We can therefore take $`Q(x)=f(x)+G(x)`$ where $`G(x)=_0^xg(s)𝑑s`$. We will use condition (ii) on $`q`$ to show that $`G`$ is bounded. Since $`\left[f^{}\right]=0`$ for any $`fL^2()`$, we have $`\left[q\right]=\left[g\right]`$ and $`\left[g\right]`$ exists. Since, also, $`gL^2()`$, the existence of $`\left[g\right]`$ implies that $`g`$ is conditionally integrable. Thus, $`lim_n\mathrm{}\alpha _n`$ exists where $`\alpha _n:=_n^ng(x)𝑑x`$. This is equivalent to the existence of $`lim_{a+\mathrm{}}_a^ag(x)𝑑x`$ if $`gL^2()`$. We need to show that the numbers $`\alpha _n^+=_0^ng(x)𝑑x`$ and $`\alpha _n^{}=_n^0g(x)𝑑x`$ are also bounded. Let $`\left\{\eta _n\right\}_{n1}`$ be a sequence of $`𝒞_0^{\mathrm{}}()`$ functions with $`0\eta _n(x)1`$, $`\eta _n(x)=1`$ for $`x[0,n]`$, $`\eta _n(x)=0`$ for $`x\backslash [1,n+1]`$, and $`\left|\eta _n^{}(x)\right|2`$. Since $`(L_q\eta _n,\eta _n)0`$,
$$(f,2\eta _n\eta _n^{})+g\eta _n^2𝑑x\eta _n^{}_{L^2()}^2.$$
Since $`\eta _n^{}_{L^2()}4`$ and
$$\left|(f,2\eta _n\eta _n^{})\right|4_1^0\left|f(x)\right|𝑑x+4_n^{n+1}\left|f(x)\right|𝑑x8f_{L^2()}$$
as well as
$$g\eta _n^2𝑑x=\alpha _n^++_1^0g(x)\eta _n^2(x)𝑑x+_n^{n+1}g(x)\eta _n^2(x)𝑑x\alpha _n^++2g_{L^2()}$$
it then follows that $`\alpha _n^+C`$ with $`C`$ independent of $`n`$. A similar argument shows that $`\alpha _n^{}C`$ with $`C`$ independent of $`n`$. As $`\alpha _n=\alpha _n^++\alpha _n^{}`$ we conclude that the sequences $`\left\{\alpha _n^+\right\}`$ and $`\left\{\alpha _n^{}\right\}`$ are both bounded, so $`G`$ is bounded. Since $`L_q0`$ we have $`r^{}+r^2=q`$ for some $`rL_{\mathrm{loc}}^2()`$ by Theorem 1.1, and applying Lemma 4.1 we conclude that $`rL^2()`$. Hence $`q\mathrm{Im}(B_0)`$.
On the other hand, if $`q\mathrm{Im}(B_0)`$, then $`L_q0`$ by Theorem 1.1 and $`q=r^{}+r^2`$ for some $`rL^2(),`$ hence by properties (c) and (d) of the special integral, $`\left[q\right]`$ exists and $`\left[q\right]=r_{L^2()}^20`$. ∎
###### Corollary 4.6.
An odd distribution $`qH^1()`$ cannot be in $`\mathrm{Im}(B_0)`$ unless $`q0.`$
###### Remark 4.7.
Generally, the condition $`L_q0`$ can be considered as a weak form of positivity for $`q`$. If it is satisfied then the existence of the special integral $`[q]`$ for $`qH^1()`$ implies much stronger existence-of-limit type results. For example, let us take any family of functions $`\chi _{T_1,T_2}𝒞_0^{\mathrm{}}()`$, $`T_1,T_2`$, such that $`\chi _{T_1,T_2}=1`$ on $`[T_1,T_2]`$, $`\chi _{T_1,T_2}=0`$ on $`(T_11,T_2+1)`$, and the derivatives $`\chi _{T_1,T_2}^{}(x)`$ are uniformly bounded. Then $`Ł_q0`$ and existence of $`[q]`$ imply the existence of the limit
$$\underset{T_1,T_2+\mathrm{}}{lim}q,\chi _{T_1,T_2},$$
which in case $`qL^2()`$ is equivalent to the existence of the limit
$$\underset{T_1,T_2+\mathrm{}}{lim}_{T_1}^{T_2}q(x)𝑑x.$$
To prove the above statements we can, for example, use Theorem 4.5 to find
$`rL^2()`$, such that $`q=r^{}+r^2`$, and the result easily follows.
We now consider the restriction $`B_\beta :H^\beta ()H^{\beta 1}()`$ for $`\beta >0`$.
###### Lemma 4.8.
Let $`\beta 0`$. If $`q\mathrm{Im}(B_0)H^{\beta 1}()`$ and $`rL_{\mathrm{loc}}^2()`$ is a solution of the Riccati equation $`r^{}+r^2=q`$, then $`rH^\beta ()`$.
###### Proof.
By Corollary 4.2, the result holds for $`\beta =0`$. Hence, it suffices to prove that in case the claimed result holds for a given $`\beta _00`$, it also holds for any $`\beta [\beta _0,\beta _0+1/4]`$. So, assume that $`q:=r^{}+r^2H^{\beta 1}()`$ with $`\beta _0<\beta \beta _0+\frac{1}{4}`$ and $`rH^{\beta _0}()`$. Then $`r^2`$ belongs to $`H^{1/2\delta }()`$, $`H^{2\beta _01/2}()`$, $`H^{1/2\delta }()`$, or $`H^{\beta _0}()`$ respectively when $`\beta _0=0`$, $`0<\beta _0<1/2`$, $`\beta _0=1/2`$, or $`\beta _0>1/2`$ (see (2.1)-(2.4)). In the first and last cases, $`\delta >0`$ can be chosen arbitrarily small. Then $`r^{}=qr^2`$ is in $`H^{s1}()`$ with $`s=\mathrm{min}(\beta ,1/2\delta ),\mathrm{min}(\beta ,2\beta _0+1/2)`$, $`\mathrm{min}(\beta ,3/2\delta )`$ or $`\mathrm{min}(\beta ,\beta _0+1)`$ respectively. As $`rL^2()`$ it then follows that $`rH^s()`$ and since $`\beta _0\beta \beta _0+1/4`$ implies $`s=\beta `$ in all cases, we get the desired result. ∎
*Proof of Theorem 1.2.* First, suppose that $`qH^{\beta 1}()`$ satisfies conditions (i) and (ii) of Theorem 1.2. From the trivial inclusion $`H^\beta ()H^0()`$ and Proposition 4.3, it follows that $`q\mathrm{Im}(B_0)`$, i.e., $`q=r^{}+r^2`$ for some function $`rL^2()`$. Applying Lemma 4.8 we see that $`rH^\beta ()`$, so $`q\mathrm{Im}(B_\beta )`$ as claimed. Second, if $`q\mathrm{Im}(B_\beta )`$, then $`L_q0`$ by Theorem 1.1 and $`q=f^{}+g`$ with $`f=rL^2()`$ and $`g=r^2L^1()`$. $`\mathrm{}`$
## 5. Geometry of the Miura Map
In this section, we prove Theorem 1.3. According to Proposition 3.2, $`B^1(q)`$ is homeomorphic to the set of positive solutions $`y`$ of $`L_qy=0`$ with $`yH_{\mathrm{loc}}^1()`$ and $`y(0)=1`$. As before we denote this set by $`Pos(q)`$. We will show that $`Pos(q)`$ is either a point or homeomorphic to a line segment.
Suppose that $`Pos(q)`$ is nonempty and choose $`y_1Pos(q)`$. Using the Wronskian we can find another solution
$$y_2(x)=y_1(x)_0^xy_1(s)^2𝑑s.$$
The general solution to $`L_qy=0`$ is then
$$y(x)=y_1(x)\left(c_1+c_2F(x)\right)$$
where
$$F(x)=_0^xy_1(s)^2𝑑s.$$
Observe that $`F`$ is a monotone increasing function with $`F(0)=0`$. If we define numbers $`m_\pm (0,+\mathrm{}]`$ by
(5.1)
$$m_+=\underset{x+\mathrm{}}{lim}F(x)$$
and
(5.2)
$$m_{}=\underset{x\mathrm{}}{lim}F(x),$$
then $`F`$ takes values in $`(m_{},m_+)`$. We will set $`m_+^1=0`$ if $`m_+=+\mathrm{}`$, and similarly for $`m_{}^1`$. The conditions $`y(0)=1`$ and $`y(x)>0`$ for all $`x`$ determine that any $`yPos(q)`$ is written
$$y(x)=y_1(x)\left(1+cF(x)\right)$$
with $`c[m_+^1,m_{}^1]`$. Letting
(5.3) $`y_+(x)`$ $`=y_1(x)\left(1m_+^1F(x)\right)`$
(5.4) $`y_{}(x)`$ $`=y_1(x)\left(1+m_{}^1F(x)\right)`$
we see that
(5.5) $`y_+(x)`$ $`y(x)y_{}(x),x>0`$
(5.6) $`y_{}(x)`$ $`y(x)y_+(x),x<0`$
for any $`yPos(q)`$, and
$$Pos(q)=\{\theta y_++(1\theta )y_{}:\theta [0,1]\}.$$
Thus, either (i) $`m_+=m_{}=+\mathrm{}`$, $`y_+=y_{}`$ and $`Pos(q)`$ consists of a single element, or (ii) at least one of $`m_\pm `$ is finite, $`y_+y_{}.`$ Noting that $`\theta \theta y_++(1\theta )y_{}`$ is a continuous map from $`[0,1]`$ to the Hausdorff space $`H_{\mathrm{loc}}^1()`$ we see that $`Pos(q)`$ is homeomorphic to the interval $`[0,1]`$. We have proved:
###### Lemma 5.1.
Suppose that $`qH_{\mathrm{loc}}^1()`$ is a real-valued distribution and $`L_q0`$. Then $`Pos(q)`$ is either a point or homeomorphic to a line segment.
Next, we show that the sets $`E_1`$ and $`E_2`$ defined in (1.5) and (1.6) are both dense in $`H_{\mathrm{loc}}^1()`$. We begin with a simple lemma which will be useful in the proof of Theorem 1.3.
###### Lemma 5.2.
There exists a family of potentials $`\left\{w_\epsilon \right\}_{\epsilon >0}`$ contained in $`𝒞_0^{\mathrm{}}()E_2`$ so that (i) $`supp(w_\epsilon )[\epsilon ^1,\epsilon ^1]`$ and (ii) $`w_\epsilon _{H^\beta ()}0`$ as $`\epsilon 0`$ for any $`\beta `$.
###### Proof.
Let $`y𝒞^{\mathrm{}}()`$ with $`y(x)=1`$ for $`x<1`$, $`y(x)=x`$ for $`x>1`$, and $`y(x)>0`$ for any $`x`$. The potential $`w(x)=y^{\prime \prime }(x)/y(x)`$ has $`y`$ as a positive solution of $`L_wy=0`$ and $`w=B(r)`$ with $`r=y^{}/y`$. Since $`_0^{\mathrm{}}y(s)^2𝑑s<\mathrm{}`$, it follows from the remarks preceding Lemma 5.1 that $`wE_2`$. Now let $`y_\epsilon (x)=y\left(\epsilon x\right)`$ and $`w_\epsilon (x)=y_\epsilon ^{\prime \prime }(x)/y_\epsilon (x)`$. Then $`w_\epsilon E_2`$ with support in $`[\epsilon ^1,\epsilon ^1]`$, proving (i). To prove (ii), note that $`w_\epsilon (x)=\epsilon ^2w\left(\epsilon x\right)`$ so that for any nonnegative integer $`j`$,
$$_x^jw_\epsilon _{L^2()}^2=\epsilon ^{3+2j}_x^jw_{L^2()}^2.$$
Since $`u_{H^\alpha ()}u_{H^\beta ()}`$ for $`\alpha <\beta `$ and $`uH^\beta ()`$, this shows that $`w_\epsilon _{H^\beta ()}0`$ as $`\epsilon 0,`$ for any $`\beta `$. ∎
###### Lemma 5.3.
$`B(𝒞_0^{\mathrm{}}())`$ is dense in $`\mathrm{Im}(B)`$ and $`B(𝒞_0^{\mathrm{}}())E_1`$.
###### Proof.
Let $`q\mathrm{Im}(B)`$ and let $`rB^1(q)`$. Let $`\left\{r_n\right\}𝒞_0^{\mathrm{}}()`$ with $`r_nr`$ in $`L_{\mathrm{loc}}^2()`$. By the continuity of the Miura map, $`B(r_n)B(r)`$ in $`H_{\mathrm{loc}}^1()`$. Thus $`B\left(𝒞_0^{\mathrm{}}()\right)`$ is dense in $`\mathrm{Im}(B)`$. If $`q=B(r)`$ for $`r𝒞_0^{\mathrm{}}()`$, then $`Pos(q)`$ contains the element $`y_1(x)=\mathrm{exp}\left(_0^xr(s)𝑑s\right)`$ which is bounded above and below by strictly positive constants. It follows that $`m_+=m_{}=+\mathrm{}`$ (see (5.1) and (5.2)), By the analysis of positive solutions preceding Lemma 5.1, $`y_+=`$ $`y_{}`$ and $`Pos(q)`$ consists of a single point. Hence $`B(𝒞_0^{\mathrm{}}())E_1`$. ∎
On the other hand:
###### Lemma 5.4.
$`E_2`$ is dense in $`\mathrm{Im}(B)`$.
###### Proof.
Since $`B(𝒞_0^{\mathrm{}}())`$ is dense in $`\mathrm{Im}(B)`$, it suffices to show that for any $`qB(𝒞_0^{\mathrm{}}())`$ there is a sequence of elements $`q_n`$ from $`E_2`$ with $`q_nq`$ in $`H_{\mathrm{loc}}^1()`$ as $`n\mathrm{}`$. Suppose that $`qB(𝒞_0^{\mathrm{}}())`$ with support in $`[a,a]`$ for $`a>0`$ and consider the sequence
$$q_n=q+v_n$$
for $`n1`$, where
$$v_n(x)=w_{1/n}(xa2n)$$
and $`w_\epsilon `$ is the family constructed in Lemma 5.2. Then $`v_n0`$ for $`n\mathrm{}`$ in $`H^\beta ()`$ for any $`\beta .`$ Let $`pH_{\mathrm{loc}}^1()`$ be the unique positive solution to $`L_qp=0`$ with $`p(0)=1`$; note that $`p(x)`$ is constant away from the support of $`q`$. If $`y_\epsilon `$ is the positive solution for $`w_\epsilon `$ constructed in Lemma 5.2, it is easily seen that the function
(5.7)
$$z_n(x)=\{\begin{array}{cc}p(x),\hfill & x<a+1,\\ p(a+1)y_{1/n}\left(xa2n\right),\hfill & xa+1,\end{array}$$
is a positive solution to $`L_{q_n}y=0`$. It follows from (5.7) and the fact that $`y_{1/n}(x)=x/n`$ for $`x`$ large and positive that $`_0^{\mathrm{}}z_n(s)^2𝑑s<\mathrm{}`$. Thus, by the analysis of positive solutions preceding Lemma 5.1, $`q_nE_2`$. Since $`v_n0`$ in $`H^\beta ()`$ for any $`\beta `$, $`q_nq0`$ in $`H_{\mathrm{loc}}^1()`$. ∎
*Proof of Theorem 1.3.* That $`\mathrm{Im}(B)=E_1E_2`$ follows from Lemma 5.1 and Proposition 3.2. The density statements were proved in Lemmas 5.3 and 5.4. $`\mathrm{}`$
We close this section with some further remarks on the dichotomy of the Miura map. First, we give a version of Theorem 1.3 for the restriction of the Miura map to $`H^\beta ()`$, $`\beta 0`$.
###### Theorem 5.5.
$`\mathrm{Im}(B_\beta )=E_{1,\beta }E_{2,\beta }`$ where $`E_{j,\beta }=E_j\mathrm{Im}(B_0)H^{\beta 1}(),`$ $`j=1,2`$. Moreover $`E_{j,\beta }`$ is dense in $`\mathrm{Im}(B_\beta )`$ for $`j=1,2`$.
###### Proof.
The first statement follows from the fact, established in Theorem 1.2, that $`\mathrm{Im}(B_\beta )=\mathrm{Im}(B_0)H^{\beta 1}()`$. The proofs of Lemmas 5.3 and 5.4 can be adapted with trivial changes to show the density of $`E_{j,\beta }`$ in $`\mathrm{Im}(B_\beta )`$. ∎
Finally, let
(5.8)
$$\lambda _0(q)=inf\{(L_q\phi ,\phi ):\phi 𝒞_0^{\mathrm{}}(),\phi _{L^2()}=1\},$$
or, equivalently,
(5.9)
$$\lambda _0(q)=inf\{\frac{𝔱_q(\psi ,\psi )}{(\psi ,\psi )}:\psi H_{\mathrm{comp}}^1()\{0\}\},$$
and define the sets
$$E_{}=\{q\mathrm{Im}(B):\lambda _0(q)=0\}$$
and
$$E_>=\{q\mathrm{Im}(B):\lambda _0(q)>0\}.$$
If $`q`$ has compact support, it is clear that $`\lambda _0(q)=0`$ since we can choose test functions whose support is disjoint from the support of $`q`$ and
$$inf\{\phi ^{}^2:\phi 𝒞_0^{\mathrm{}}(),\phi _{L^2()}=1\}=0.$$
Note that the map $`E_{}\times _{>0}E_>`$ given by
$$(q,c)q+c$$
is a continuous, bijective map onto $`E_>`$.
###### Theorem 5.6.
(i) $`E_>E_2`$ and $`E_1E_{}`$.
(ii) $`E_2E_{}\mathrm{}`$, i.e., $`E_{}`$ is not a fold of the dichotomy. Moreover, $`E_{}`$ is dense in $`\mathrm{Im}(B)`$, and $`E_1`$ and $`E_2E_{}`$ are dense in $`E_{}`$.
###### Remark 5.7.
The fact that $`E_2E_{}\mathrm{}`$ has already been observed by Murata , Remark to Theorem 2.2, in his investigation of critical and subcritical potentials – see Appendix C.
The proof of Theorem 5.6 will rely on the following proposition.
###### Proposition 5.8.
Suppose that $`qH_{\mathrm{loc}}^1()`$ and $`\lambda _0(q)>0`$. Then the equation $`L_qy=0`$ has two linearly independent positive solutions $`y_1,y_2H_{\mathrm{loc}}^1()`$.
###### Remark 5.9.
For potentials $`qL_{\mathrm{loc}}^1()`$, the result above is due to Murata , Remark after Theorem 2.7.
In the proof of Proposition 5.8 we will use
###### Lemma 5.10.
Assume that $`yH_{\mathrm{loc}}^1()`$, and there exists a discrete subset $`S`$, such that $`L_qy=0`$ on $`S`$. Then
(5.10)
$$L_qy=\underset{zS}{}(u(z0)u(z+0))\delta (z),$$
where $`u=y^{}Qy`$, $`Q^{}=q`$ as in (2.7). In other words,
(5.11)
$$(y,L_q\phi )=\underset{zS}{}(u(z0)u(z+0))\overline{\phi (z)},$$
for every $`\phi 𝒞_0^{\mathrm{}}()`$.
###### Proof.
Using partition of unity, we can split any function $`\phi 𝒞_0^{\mathrm{}}()`$ into a finite sum of functions $`\phi _k`$ such that for every $`k`$ a neighborhood of $`\mathrm{supp}\phi _k`$ contains at most one point from $`S`$. Therefore, taking into account translation invariance, we see that it suffices to consider the case when $`S=\{0\}`$. So we will assume that $`yH_{\mathrm{loc}}^1()`$ and $`L_qy=0`$ on $`0`$.
Integrating by parts (see (2.5)) and using (2.7), we get
$`(y,L_q\phi )=(y^{},\phi ^{})+(qy,\phi )`$
$`=(u+Qy,\phi ^{})+(qy,\phi )=(u,\phi ^{})+(qy(Qy)^{},\phi )`$
$`=(u,\phi ^{})(Qy^{},\phi )=(u,\phi ^{})(Q^2y+Qu,\phi ).`$
Integrating by parts in the first term in the right hand side we obtain
$$(u,\phi ^{})=_{\mathrm{}}^0u\overline{\phi ^{}}𝑑x+_0^{\mathrm{}}u\overline{\phi ^{}}𝑑x=(u(0)u(+0))\overline{\phi (0)}([u^{}],\phi ),$$
where $`[u^{}]`$ is the locally integrable function on $``$ which coincides with $`u^{}`$ on $`\{0\}`$. Since $`[u^{}]=Q^2yQu`$ due to (2.7), we finally obtain
$$(y,L_q\phi )=(u(0)u(+0))\overline{\phi (0)},$$
as required. ∎
###### Corollary 5.11.
Let $`y`$ satisfy the conditions of Lemma 5.10 and have a compact support (so that $`S`$ can be taken finite). Then
(5.12)
$$𝔱_q(y,y)=_{}(|y^{}|^2+q|y|^2)𝑑x=\underset{zS}{}(u(z0)u(z+0))\overline{y(z)}.$$
###### Proof.
Taking limit in (5.11) over a sequence $`\phi _k`$ converging to $`y`$ in $`H_{\mathrm{comp}}^1()`$, we obtain (5.12). ∎
###### Proof of Proposition 5.8.
Using notations from the proof of Proposition 3.5 (see (3.4)), for any $`c>0`$ define a test function (to use in (5.9))
(5.13)
$$\psi _c(x)=\{\begin{array}{cc}y_c(x),\hfill & x(c,0);\hfill \\ y_c(x),\hfill & x[0,c);\hfill \\ 0,\hfill & x(c,c).\hfill \end{array}$$
Clearly, $`\psi _cH_{\mathrm{comp}}^1()`$, $`\psi _c(0)=1`$, and $`L_q\psi _c=0`$ on $`S`$ where $`S=\{c,0,c\}`$.
Applying Corollary 3.4 to $`y_cy_c^{}`$, we see that $`c\psi _c(x)`$ is an increasing function of $`c>0`$ for any fixed $`x`$. Therefore, the $`L^2`$-norm $`\psi _c`$ increases with $`c`$ as well.
By Corollary 5.11 we have
$$𝔱_q(\psi _c,\psi _c)=u_c(0)u_c(0).$$
It follows from Lemma 2.5 that $`u_c(0)`$ decreases and $`u_c(0)`$ increases as $`c`$ increases. Therefore, $`c𝔱_q(\psi _c,\psi _c)`$ is decreasing as $`c`$ increases. It follows that the fraction in (5.9) is decreasing as well. To prove the desired statement, it is enough to establish that the limit of this fraction is $`0`$ as $`c+\mathrm{}`$, provided we know that the equation $`L_qy=0`$ has only one positive solution with $`y(0)=1`$. To this end note that the limits of $`y_c`$ and $`y_c`$ exists and are both positive solutions, according to the arguments given in the proof of Proposition 3.5. Due to our uniqueness of positive solution assumption these limits should coincide. But then we should also have
$$\underset{c+\mathrm{}}{lim}u_c(0)=\underset{c+\mathrm{}}{lim}u_c(0),$$
because $`y_c(0)=y_c(0)=1`$ and the map $`y\{y(0),u(0)\}`$ is a linear topological isomorphism between the space of all solutions of $`L_qy=0`$ with the $`H_{\mathrm{loc}}^1()`$-topology and the space $`^2`$ It follows that
$$\underset{c+\mathrm{}}{lim}𝔱_q(\psi _c,\psi _c)=0,$$
which implies the desired statement. ∎
###### Proof of Theorem 5.6.
To prove part (i), it is enough to show that $`E_>E_2`$ since it then follows by taking complements that $`E_1E_{}`$. In Proposition 5.8, we established that any $`q\mathrm{Im}(B)`$ with $`\lambda _0(q)>0`$ has two linearly independent, positive solutions of $`L_qy=0`$ in $`H_{\mathrm{loc}}^1()`$, so $`E_>E_2`$.
To prove part (ii), we first note that, by Lemma 5.2, there are compactly supported potentials in $`E_2`$, and by the remark above, $`\lambda _0(q)=0`$ for such potentials, so $`E_2E_{}`$ is nonempty. Next, note that $`B\left(𝒞_0^{\mathrm{}}()\right)𝒞_0^{\mathrm{}}()`$ so $`B\left(𝒞_0^{\mathrm{}}()\right)E_{}`$. On the other hand, by Lemma 5.3, $`B\left(𝒞_0^{\mathrm{}}()\right)`$ is dense in $`\mathrm{Im}(B)`$, so $`E_{}`$ is dense in $`\mathrm{Im}(B)`$. We have already shown that $`E_1`$ is dense in $`\mathrm{Im}(B)`$, so $`E_1`$ is also dense in $`E_{}`$ by part (i). The proof of Lemma 5.4 shows that $`E_2E_{}`$ is dense in $`E_{}.`$
###### Remark 5.12.
Note that the map $`\mathrm{\Phi }:E_{}\times _0\mathrm{Im}(B)`$ defined by $`(q,c)q+c`$ is continuous and bijective, but not a homeomorphism. Otherwise, $`\mathrm{\Phi }(E_{}\times \left\{0\right\})\mathrm{Im}(B)`$ would be closed, and, as $`E_1`$ is dense in $`\mathrm{Im}(B)`$, we conclude that $`E_{}=\mathrm{Im}(B)`$, a contradiction. The interpretation that $`E_2`$ is at least “one dimension larger” than $`E_1`$ could therefore be somewhat misleading. Note that the inverse of $`\mathrm{\Phi }`$ is given by $`\mathrm{\Phi }^1:\mathrm{Im}(B)E_{}\times _0`$, $`q(q\lambda _0(q),\lambda _0(q))`$. Hence, $`\mathrm{\Phi }^1`$ not being continuous means that $`q\lambda _0(q)`$ is not continuous in $`H_{\mathrm{loc}}^1()`$.
## 6. Application to KdV
In this section we apply our results on the Miura map to prove existence of solutions of the Korteweg-de Vries equation in $`H^1()`$ for initial data in the range $`\mathrm{Im}(B_0)`$ of the Miura map $`B_0:L^2()H^1()`$. We follow the approach of Tsutsumi , who proved such an existence result for initial data a positive, finite Radon measure on $``$. His arguments combined with our results on the Miura map $`B_0`$ lead to the following theorem. Recall that, for a real-valued distribution $`uH^1()`$, $`\left[u\right]`$ denotes the special integral of $`u`$ (see Theorem 4.5 and the discussion that precedes it).
###### Theorem 6.1.
Assume that $`u_0\mathrm{Im}(B_0)`$. Then there exists a global weak solution of KdV with $`u(t)\mathrm{Im}\left(B_0\right)`$ for all $`t`$. More precisely:
(i) $`uL^{\mathrm{}}(,H^1())L_{\mathrm{loc}}^2(^2)`$,
(ii) for all functions $`\phi 𝒞_0^{\mathrm{}}(^2)`$, the identity
$$_{}_{}\left(u\phi _tu\phi _{xxx}+3u^2\phi _x\right)𝑑x𝑑t=0$$
holds,
(iii) $`lim_{t0}u(t)=u_0`$ in $`H^1()`$, and
(iv) $`0\left[u(t)\right]\left[u_0\right]`$ for all $`t`$ and $`lim_{t0}\left[u(t)\right]=\left[u_0\right].`$
###### Remark 6.2.
Recall that $`u_0\mathrm{Im}(B_0)`$ means that $`u_0H^1()`$ with the property that $`L_{u_0}0`$ and $`\left[u_0\right]`$ exists. Instead of the assumption for $`\left[u_0\right]`$ to exist, one can equivalently assume that $`u_0=f^{}+g`$ for some functions $`fL^2()`$ and $`gL^1()`$ – see Theorem 1.2 and Proposition 4.3.
To prove Theorem 6.1, we need to recall a result of Kato and, independently, of Kruzhkov and Faminskiĭ – see also and . Consider the modified Korteweg-de Vries equation (mKdV)
(6.1)
$$_tv=_x^3v+6v^2_xv$$
with initial data
(6.2)
$$v(0)=v_0.$$
###### Theorem 6.3.
, Let $`v_0L^2()`$. Then the initial value problem (6.1)-(6.2) has a weak solution in $`L^{\mathrm{}}(,L^2())`$. More precisely, $`v`$ satisfies:
(i) $`vL^{\mathrm{}}(,L^2())L_{\mathrm{loc}}^2(,H_{\mathrm{loc}}^1())`$,
(ii) the identity
$$_{}_{}\left(v\phi _tv\phi _{xxx}+2v^3\phi _x\right)𝑑t𝑑x=0$$
holds for all $`\phi 𝒞_0^{\mathrm{}}(^2)`$,
(iii) $`lim_{t0}v(t)=v_0`$ in $`L^2()`$, and
(iv) $`v(t)_{L^2()}v_0_{L^2()}`$ for all $`t`$.
The following result improves the one of Tsutsumi by adapting it to our more general setting, and relies on the identity (1.2).
###### Proposition 6.4.
Let $`v=v(t)`$ be a solution of (6.1)-(6.2) with $`v_0L^2()`$ and the properties listed in Theorem 6.3. Let $`u_0:=v_0^{}+v_0^2`$. Then $`u:=v^{}+v^2`$ is a solution of KdV. More precisely:
(i) $`uL^{\mathrm{}}(,H^1())L_{\mathrm{loc}}^2(^2)`$,
(ii) the identity
$$_{}_{}\left(u\phi _tu\phi _{xxx}+3u^2\phi _x\right)𝑑t𝑑x=0$$
holds for all $`\phi 𝒞_0^{\mathrm{}}(^2)`$,
(iii) $`lim_{t0}u(t)=u_0`$ in $`H^1()`$,
(iv) $`0\left[u(t)\right]\left[u_0\right]`$ for all $`t,`$ and $`lim_{t0}\left[u(t)\right]=\left[u_0\right].`$
###### Proof.
Statement (i) follows from Theorem 6.3(i) together with the fact that $`B_0:L^2()H^1()`$ is a bounded, continuous map – see Proposition 2.3. Statement (iii) follows from Theorem 6.3(iii) and the continuity of $`B_0`$ whereas the claimed inequality in (iv) follows from Theorem 6.3(iv) and the fact that $`\left[u(t)\right]=v(t)_{L^2()}^2`$ – see formula (4.4). To prove the second statement in (iv), note that by Theorem 6.3(iii), $`lim_{t0}v(t)_{L^2()}=v_0_{L^2()}`$. As $`\left[u(t)\right]=v(t)_{L^2()}^2,`$ the second statement in (iv) then follows as well. Statement (ii) is proved in . For the convenience of the reader we include a detailed proof of it.
Let $`\rho :\times `$ be the smooth mollifier, i.e. a smooth positive function with support in the unit disc in $`^2`$, $`\rho (0,0)>0,`$ and normalized by $`_^2\rho (t,x)𝑑t𝑑x=1`$. For $`\epsilon >0`$, set
$$\rho _\epsilon (t,x):=\frac{1}{\epsilon ^2}\rho (\frac{t}{\epsilon },\frac{x}{\epsilon }).$$
Given the solution $`v(t)`$ of mKdV, define for $`(t,x)^2`$
$`v_\epsilon (t,x)`$ $`:=\left(\rho _\epsilon v\right)(t,x)`$
$`={\displaystyle _^2}\rho (ts,xy)v(s,y)𝑑s𝑑y.`$
Note that for any $`\epsilon >0`$ and $`(t,x)^2`$, $`\rho _\epsilon (ts,xy)𝒞_0^{\mathrm{}}(^2)`$ as a function of $`(s,y)^2`$. Further define the function $`u_\epsilon 𝒞_0^{\mathrm{}}(^2),`$
$$u_\epsilon :=\frac{}{x}v_\epsilon +v_\epsilon ^2.$$
According to (1.2), one has
(6.3) $`{\displaystyle \frac{}{t}}u_\epsilon +{\displaystyle \frac{^3}{x^3}}u_\epsilon 6u_\epsilon {\displaystyle \frac{}{x}}u_\epsilon `$ $`=\left({\displaystyle \frac{}{x}}+2v_\epsilon \right)\left({\displaystyle \frac{}{t}}v_\epsilon +{\displaystyle \frac{^3}{x^3}}v_\epsilon 6\rho _\epsilon \left(v^2{\displaystyle \frac{}{x}}v\right)\right)`$
$`+6\left({\displaystyle \frac{}{x}}+2v_\epsilon \right)\left(\rho _\epsilon \left(v^2{\displaystyle \frac{}{x}}v\right)v_\epsilon ^2{\displaystyle \frac{}{x}}v_\epsilon \right).`$
By assumption, $`v`$ is a weak solution of mKdV, hence
$`{\displaystyle \frac{}{t}}v_\epsilon +{\displaystyle \frac{^3}{x^3}}v_\epsilon 6\rho _\epsilon \left(v^2{\displaystyle \frac{}{x}}v\right)`$ $`=\rho _\epsilon \left({\displaystyle \frac{}{t}}v+{\displaystyle \frac{^3}{x^3}}v6v^2{\displaystyle \frac{}{x}}v\right)`$
$`=0.`$
Multiplying (6.3) by an arbitrary test function $`\phi 𝒞_0^{\mathrm{}}(^2)`$ and integrating by parts, one obtains
(6.4) $`{\displaystyle _^2}\left(u_\epsilon {\displaystyle \frac{}{t}}\phi u_\epsilon {\displaystyle \frac{^3}{x^3}}\phi +3u_\epsilon ^2{\displaystyle \frac{}{x}}\phi \right)𝑑t𝑑x`$
$`=6{\displaystyle _^2}\left[\rho _\epsilon \left(v^2{\displaystyle \frac{}{x}}v\right)v_\epsilon ^2{\displaystyle \frac{}{x}}v_\epsilon \right]\left({\displaystyle \frac{}{x}}\phi +2v_\epsilon \phi \right)𝑑t𝑑x`$
$`=2{\displaystyle _^2}\left(\rho _\epsilon v^3v_\epsilon ^3\right)\left({\displaystyle \frac{^2}{x^2}}\phi 2\phi {\displaystyle \frac{}{x}}v_\epsilon 2v_\epsilon {\displaystyle \frac{}{x}}\phi \right)𝑑t𝑑x.`$
By Theorem 6.3(i),
$$\frac{}{x}v_\epsilon \frac{}{x}v\text{in }L_{\mathrm{loc}}^2(^2).$$
Lemma 6.5 below together with Theorem 6.3(i) implies that $`vL_{\mathrm{loc}}^6(^2)`$. Hence,
$$v_\epsilon ^2v^2,v_\epsilon ^3v^3\text{in }L_{\mathrm{loc}}^2(^2).$$
Combining all of this, one obtains
$`\rho _\epsilon v^3v_\epsilon ^3`$ $`=\left(\rho _\epsilon v^3v^3\right)\left(v_\epsilon ^3v^3\right)0\text{ as }\epsilon 0\text{ in }L_{\mathrm{loc}}^2(^2)`$
and
$$2\phi \frac{}{x}v_\epsilon +2v_\epsilon \frac{}{x}\phi 2\phi \frac{}{x}v+2v\frac{}{x}\phi \text{ in }L^2(^2).$$
Since $`u_\epsilon =_xv_\epsilon +v_\epsilon ^2`$, it follows that
$$u_\epsilon u\text{ in }L_{\mathrm{loc}}^2(^2)$$
and, as $`supp\phi `$ is compact,
$$_^2\left(\rho _\epsilon v^3v_\epsilon ^3\right)\left(\frac{^2}{x^2}\phi 2\phi \frac{}{x}v_\epsilon 2v_\epsilon \frac{}{x}\phi \right)𝑑t𝑑x0$$
as $`\epsilon 0`$. Therefore, taking the limit $`\epsilon 0`$ in (6.4), we conclude that
$$_^2\left(u\frac{}{t}\phi u\frac{^3}{x^3}\phi +3u^2\frac{}{x}\phi \right)𝑑t𝑑x=0.$$
###### Lemma 6.5.
Let $`Q:=I\times J`$ with $`I:=[T,T]`$ and $`J:=[R,R]`$ where $`T>0`$ and $`R>0`$. Let
$$_{I,J}=L^2(I,H^1(J))L^{\mathrm{}}(I,L^2(J))$$
with norm
$$ess\underset{tI}{sup}f(t,)_{L^2(J)}+\left(_If(t,)_{H^1(J)}^2𝑑t\right)^{1/2}.$$
Then, for any $`f_{I,J}`$, the inequality
(6.5)
$$f_{L^6\left(Q\right)}^6Cess\underset{tI}{sup}f(t,)_{L^2(J)}^4_If(t,)_{H^1(J)}^2𝑑t$$
holds, where $`C>0`$ is a constant which depends only on $`R`$. In particular, $`_{I,J}`$ embeds continuously into $`L^6(Q)`$.
###### Proof.
Let us assume first that $`f𝒞_0^{\mathrm{}}(^2)`$. By the standard Sobolev embedding theorem, the space $`H^{1/3}(J)`$ is densely and continuously embedded in $`L^6(J)`$ (see e.g. , Theorem 1, p 82),
$$g_{L^6(J)}C_J^{}g_{H^{1/3}(J)}$$
for all $`gH^{1/3}(J)`$ and some positive constant $`C_J^{}`$. Further, by interpolation, one has for any $`gH^1(J)`$ (see e.g. , Remark 2, p 87)
$$g_{H^{1/3}(J)}C_J^{\prime \prime }g_{H^1(J)}^{1/3}g_{L^2(J)}^{2/3}.$$
Setting $`C=\left(C_J^{}C_J^{\prime \prime }\right)^6`$, we see that
$`{\displaystyle _{I\times J}}\left|f(t,x)\right|^6𝑑t𝑑x`$ $`{\displaystyle _I}\left(C_J^{}f(t,)_{H^{1/3}(J)}\right)^6𝑑t`$
$`C{\displaystyle _I}f(t,)_{H^1(J)}^2f(t,)_{L^2(J)}^4𝑑t.`$
By approximation, the above inequality holds for any $`f_{I,J}.`$ For such $`f`$ we have
$`{\displaystyle _I}f(t,)_{H^1(J)}^2f(t,)_{L^2(J)}^4𝑑t`$
$`ess\underset{tI}{sup}f(t,)_{L^2(J)}^4{\displaystyle _I}f(t,)_{H^1(J)}^2𝑑t.`$
Combinig the two previous inequalities ends the proof. ∎
###### Proof of Theorem 6.1.
The proof is the one given in Tsutsumi , adapted to our more general setting. By our assumption, $`B_0(v_0)=u_0`$ for some $`v_0L^2()`$. By Theorem 6.3, there exists a solution $`vL^{\mathrm{}}(,L^2())L_{\mathrm{loc}}^2(,H_{\mathrm{loc}}^1())`$ of (6.1)-(6.2). By Proposition 6.4, $`u(t):=B(v(t))`$ is a solution of KdV satisfying (i)-(iv). ∎
## Appendix A Positive Solutions for Square-Well Potentials
In this appendix we present some elementary but important examples of potentials in $`\mathrm{Im}(B)`$ and the associated positive solutions. Let
$$q_{a,b}(x)=\{\begin{array}{cc}b^2,& a<x<a\\ & \\ 0,& \left|x\right|a\end{array}$$
where $`a,b>0.`$ It is easy to see that
(A.1)
$$y_+(x)=\{\begin{array}{ccc}1/\mathrm{cosh}(ba),\hfill & & x<a\hfill \\ & & \\ \mathrm{cosh}(b(x+a))/\mathrm{cosh}(ba),\hfill & & a<x<a\hfill \\ & & \\ \left(\mathrm{cosh}(2ab)+b(xa)\mathrm{sinh}(2ab)\right)/\mathrm{cosh}(ba),\hfill & & x>a\hfill \end{array}$$
and $`y_{}(x):=y_+(x)`$ are linearly independent positive solutions of $`y^{\prime \prime }+q_{a,b}y=0`$.
If $`\lambda >0`$ and $`b=\left(\lambda /2a\right)^{1/2}`$ then $`q_{a,b}(x)𝑑x=\lambda `$. Taking $`a0`$ we recover in the limit $`q=\lambda \delta `$ where $`\delta `$ is the Dirac $`\delta `$-distribution at $`x=0`$. In this limit
$$y_+(x)=\{\begin{array}{ccc}1\hfill & & x<0\hfill \\ & & \\ 1+\lambda x\hfill & & x>0\hfill \end{array}$$
and again $`y_{}(x)=y_+(x)`$.
Let us determine the preimage of $`q=\lambda \delta `$ by the Miura map. From the explicit formulas we have $`_0^{\mathrm{}}y_+(s)^2𝑑s<\mathrm{}`$ but $`_{\mathrm{}}^0y_+(s)^2𝑑s=+\mathrm{}`$, while the reverse is true for $`y_{}`$. If $`H`$ is the Heaviside function
$$H(x)=\{\begin{array}{cc}0& x<0\\ & \\ 1& x>0\end{array}$$
then the logarithmic derivatives
$$\frac{y_+^{}(x)}{y_+(x)}=\frac{\lambda H(x)}{1+\lambda x},\frac{y_{}^{}(x)}{y_{}(x)}=\frac{y_+^{}(x)}{y_+(x)}$$
belong to $`L^2()`$. Hence
$$B^1(\lambda \delta )=\{(1\theta )\frac{\lambda H(x)}{1+\lambda x}\theta \frac{\lambda H(x)}{1\lambda x}0\theta 1\}.$$
## Appendix B Positive Schrödinger Operators
In this Appendix we provide more information about Schrödinger operators $`L_q`$ which are positive or, more generally, semibounded below, and have real potentials $`qH_{\mathrm{loc}}^1()`$. Namely, we will show that the corresponding quadratic form (defined on $`𝒞_0^{\mathrm{}}()`$) is closable and describe the domain of its closure. We will also describe the domain of the corresponding self-adjoint operator. Finally, for strictly positive $`L_q`$ (such that $`\lambda _0(q)>0`$, see (5.8), (5.9)) we construct Green’s function and use it to give an alternative proof of Proposition 5.8.
The case of semi-bounded $`L_q`$ for many purposes is reduced to the case when $`L_q0`$ or even to the case when $`L_q`$ is strictly positive (that is $`\lambda _0(q)>0`$) by adding a sufficiently large constant to $`q`$. So let us assume first that $`L_q0`$. By Theorem 1.1 there exists a function $`rL_{\mathrm{loc}}^2()`$ such that $`q=B(r)`$. Then $`L_q`$ admits a formal factorization (1.4), i.e., a presentation $`L_q=P^+P`$, where $`P=\left(_xr\right)`$ and $`P^+=\left(_x+r\right)`$, so that $`P^+`$ is the operator formally adjoint to $`P`$ in $`L^2()`$.
Clearly, $`P,P^+`$ are well defined on the space $`𝒞_0^{\mathrm{}}()`$ which is dense in $`L^2()`$, so that $`P,P^+`$ map $`𝒞_0^{\mathrm{}}()`$ to $`L^2()`$ and
$$(Pu,v)=(u,P^+v),u,v𝒞_0^{\mathrm{}}().$$
It follows that the operators $`P,P^+`$ are closable, with the closures which we will denote by $`\overline{P},\overline{P^+}`$. They also have adjoint operators in $`L^2()`$, which will be denoted $`P^{},(P^+)^{}`$. These operators are closed extensions of $`P^+,P`$ respectively. Since $`P^{},(P^+)^{}`$ are closed, we have
(B.1)
$$\overline{P}(P^+)^{},\overline{P^+}P^{}.$$
###### Lemma B.1.
(i) We have
(B.2)
$$\overline{P}=(P^+)^{},\overline{P^+}=P^{}.$$
(ii) The domains of the operators in (B.2) are as follows:
(B.3)
$$𝔇(\overline{P})=\{uL^2()W_{\mathrm{loc}}^{1,1}():PuL^2()\},$$
(B.4)
$$𝔇(\overline{P^+})=\{vL^2()W_{\mathrm{loc}}^{1,1}():P^+vL^2()\},$$
where the operators $`P,P^+`$ are applied in the usual distributional sense.
(iii) $`𝒞_0^{\mathrm{}}()`$ is an operator core for each of the operators in (B.2).
###### Proof.
It is easy to see that the right-hand sides in (B.3), (B.4) coincide with the domains of the adjoint operators $`(P^+)^{},P^{}`$ respectively. Indeed, the relation $`(P^+)^{}u=f`$ means that $`u,fL^2()`$ and for every $`\phi 𝒞_0^{\mathrm{}}()`$
$$(u,P^+\phi )=(u,_x\phi )(u,r\phi )=(f,\phi ).$$
Since $`ruL_{\mathrm{loc}}^1()`$, this is equivalent to $`_xuru=f`$, where $`_x`$ is applied in the sense of distributions. It follows that $`_xu=f+ruL_{\mathrm{loc}}^1()`$, hence $`uW_{\mathrm{loc}}^{1,1}()`$ or, equivalently, $`u`$ is absolutely continuous. This means that the right hand side of (B.3) coincides with $`𝔇((P^+)^{})`$. The same argument applies to the operator $`P^{}`$ and the right-hand side of (B.4).
Taking into account the inclusions (B.1), we see that to establish all statements of the lemma, it suffices to show that the right hand sides of (B.3), (B.4) belong to the domains of $`\overline{P},\overline{P^+}`$ respectively. This is easily done by use of Friedrichs’ mollifiers. It is essentially a special case of Friedrichs’ well-known result on equality of weak and strong extensions of differential operators, but we give the proof for the reader’s convenience. We will give the arguments for $`P`$ (the arguments for $`P^+`$ are the same).
So let us assume that
(B.5)
$$uL^2()W_{\mathrm{loc}}^{1,1}(),PuL^2().$$
We need to show that $`u`$ may be approximated by a sequence $`\left\{u_n\right\}_{n1}`$ from $`𝒞_0^{\mathrm{}}()`$ with $`u_nu`$ and $`Pu_nPu`$ in $`L^2()`$. First, we show that it suffices to consider $`u`$ satisfying (B.5) and additionally having compact support. To this end, take $`\chi 𝒞_0^{\mathrm{}}()`$. Then $`\chi u`$ also satisfies (B.5) and $`P\left(\chi u\right)=\chi ^{}u+\chi Pu`$. If $`\chi _n𝒞_0^{\mathrm{}}()`$ satisfies the conditions $`0\chi _n1`$, $`\chi _n(x)=1`$ for $`\left|x\right|n`$, $`\chi _n(x)=0`$ for $`\left|x\right|n+1`$, and $`\left|\chi _n^{}(x)\right|2`$, then $`\chi _nuu`$ in $`L^2()`$ as $`n\mathrm{}`$. Moreover, $`\chi _n^{}u0`$ and $`\chi _nPuPu`$ in $`L^2()`$ as $`n\mathrm{}`$, so $`P(\chi _nu)=\chi _n^{}u+\chi _nPuPu`$ in $`L^2()`$ as $`n\mathrm{}`$ Thus, we may assume that $`u`$ has compact support.
Given $`u`$ satisfying (B.5) and having compact support, we now use Friedrichs mollifiers to construct a sequence of approximants from $`𝒞_0^{\mathrm{}}()`$. Let $`j𝒞_0^{\mathrm{}}()`$ be a nonnegative function with $`j(x)𝑑x=1`$, and, for any $`k`$, let $`j_k(x)=kj(kx)`$, and let $`u_k=uj_k`$. Clearly $`u_k𝒞_0^{\mathrm{}}()`$ and $`u_ku`$ in $`L^2()`$. We claim that $`Pu_kPu`$ in $`L^2()`$. Since $`uW_{\mathrm{comp}}^{1,1}()`$ and $`rL_{\mathrm{loc}}^2()`$, it follows that $`ruL^1()`$. Moreover, as $`u`$ satisfies (B.5) and the support of $`u`$ is compact, $`PuL^1().`$ Therefore $`u^{}=Pu+ru`$ belongs to $`L^1()`$. Hence, $`u`$ is a bounded, continuous function which shows that $`ruL^2().`$ As a consequence, $`u^{}=Pu+ruL^2()`$. Thus $`u_k^{}u^{}`$ in $`L^2(),`$ $`u_ku`$ in $`L^{\mathrm{}}()`$ as $`k\mathrm{}`$ and so
$$Pu_kPu=\left(u_k^{}u^{}\right)+r\left(u_ku\right)$$
converges to zero in $`L^2()`$ as $`k\mathrm{}`$. ∎
Now let us recall a classical theorem of von Neumann (see also , Theorem XI.23) which asserts that if $`A`$ is a closed densely defined operator in a Hilbert space, then the (generally unbounded) operator $`H=A^{}A`$ is self-adjoint. Here the domain of $`A^{}A`$ is naturally defined as
$$𝔇(A^{}A)=\{u𝔇(A),Au𝔇(A^{})\}.$$
(Note that an essentially inverse statement also holds: if two densely defined operators $`A,A^+`$ are formally adjoint, that is
$$(Au,v)=(u,A^+v),u𝔇(A),v𝔇(A^+),$$
and $`A^+A`$ is essentially self-adjoint, then the closures $`\overline{A},\overline{A^+}`$ are adjoint to each other; see the appendix to .)
The following lemma is well-known.
###### Lemma B.2.
Let $`A`$ be a closed densely defined operator in a Hilbert space, and $`H=A^{}A`$. Denote by $`𝔱_H`$ the quadratic form of $`H`$, and let $`𝔇(𝔱_H)`$ be its domain i.e. $`𝔇(𝔱_H)=𝔇(H^{1/2})`$. Then $`𝔇(𝔱_H)=𝔇(A)`$ and
$$𝔱_H(u,u)=Au^2,u𝔇(A).$$
###### Proof.
Take the polar decomposition $`A=U|A|`$, where $`|A|=(A^{}A)^{1/2}=H^{1/2}`$, and $`U`$ partial isometry with $`\mathrm{Ker}U=\mathrm{Ker}A`$ (see e.g. Sect. VIII.9 in ). It remains to notice that $`𝔇(A)=𝔇(|A|)`$ (because $`U`$ is bounded), and
$$𝔱_H(u,u)=H^{1/2}u^2=|A|u^2=Au^2,u𝔇(A),$$
because $`U`$ is an isometry on the range of $`|A|`$. ∎
Note that a positive self-adjoint operator $`H`$ is uniquely defined by its (positive, closed) quadratic form (see e.g. Theorem VIII.15 in ).
Taking the quadratic form $`𝔱_q`$, corresponding to a potential $`qH_{\mathrm{loc}}^1()`$ and assuming that it is positive (or, more generally, semi-bounded below), we can construct a unique self-adjoint operator $`H`$ with this form. It follows from the considerations above that when the form $`𝔱_q`$ is positive, we can write this operator in the form $`H=P^{}\overline{P}`$, and the domain of $`H`$ is
(B.6)
$$𝔇(H)=\{uL^2(),(_xr)uL^2(),(_x+r)\left[(_xr)u\right]L^2()\},$$
where $`rL_{\mathrm{loc}}^2()`$ is a solution of the Riccati equation $`r^{}+r^2=q`$.
In case when $`L_q`$ is semibounded below but not positive, the arguments given above should be applied to the operator $`L_q+c`$ with $`c>0`$ such that $`L_q+c0`$, with subsequent subtracting of the same constant $`c`$ from the resulting operator (which does not change the domain of the operator). The resulting operator $`H`$ will not depend of the choice of $`c`$ because different choices of $`c`$ lead to the same (closed) quadratic form of the resulting operator.
The following lemma simplifies calculation of $`Hu`$ if we know that $`u𝔇(H)`$.
###### Lemma B.3.
Let $`qH_{\mathrm{loc}}^1()`$ be such that $`L_q0`$. Then $`H`$ can be extended to a linear operator $`\stackrel{~}{L}_q`$ with the domain
(B.7)
$$𝔇(\stackrel{~}{L}_q)=\{uH_{\mathrm{loc}}^1()L^2(),\stackrel{~}{L}_quL^2()\},$$
where $`\stackrel{~}{L}_q`$ is $`_x^2+q`$ applied in the sense of distributions, i.e. both $`_x^2`$ and $`q`$ act as linear continuous operators $`H_{\mathrm{loc}}^1()H_{\mathrm{loc}}^1()`$ ($`_x^2`$ acts as the distributional derivative, and $`q`$ acts as a multiplier in these spaces).
###### Proof.
If $`u𝔇(H)`$ (as described by (B.6)) then $`uL^2()`$ and $`f:=u^{}ruL^2()`$. Since $`ruL_{\mathrm{loc}}^1()`$, we see that $`u^{}L_{\mathrm{loc}}^1()`$, hence $`uW_{\mathrm{loc}}^{1,1}()`$, i.e. $`u`$ is absolutely continuous. But then $`ruL_{\mathrm{loc}}^2()`$, and $`uH_{\mathrm{loc}}^1()`$. Now we can conclude that
$$(_x+r)(_xr)u=(_x^2+q)u,$$
where all operations should be applied in the distributional sense. We proved that $`𝔇(H)𝔇(\stackrel{~}{L}_q)`$ and $`H`$, applied as a factorized operator (see (1.4)), is a restriction of $`\stackrel{~}{L}_q`$. ∎
###### Remark B.4.
In particular, we can apply $`H`$ on $`𝔇(H)`$ as $`_x^2+q`$ applied termwise, which is usually easier than to apply it in the factorized form. To illustrate it, note that it may easily happen that $`𝔇(\stackrel{~}{L}_q)`$ does not contain any function $`u𝒞_0^{\mathrm{}}()`$ (except $`u0`$). This is true e.g. for the potential
$$q(x)=\underset{k=1}{\overset{\mathrm{}}{}}c_k\delta (xx_k),\underset{k=1}{\overset{\mathrm{}}{}}c_k<\mathrm{},$$
where $`c_k>0`$ for all $`k`$, and the set $`\{x_k\}_{k=1}^{\mathrm{}}`$ is dense in $``$. Therefore, the same is true for $`𝔇(H)`$.
###### Remark B.5.
The operator $`\stackrel{~}{L}_q`$ may be defined on the domain (B.7) even without semi-boundednes requirement. But in this case the resulting operator (called usually “maximal operator”) in $`L^2()`$ will not necessarily be self-adjoint even if $`q𝒞^{\mathrm{}}()`$ (see e.g. Sect. X.1 in or Sect. II.1 and II.4.2 in ).
Now, note that $`\lambda _0(q)`$ (as defined in (5.8)) is the bottom of the spectrum of the self-adjoint operator $`H`$. So, if $`\lambda _0(q)>0`$, then there exists a bounded, everywhere defined linear operator $`T:=H^1`$ in $`L^2()`$. It maps $`L^2()`$ onto $`𝔇(H)`$. We will now analyze the properties of the Schwartz kernel of $`T`$, which is Green’s function for the operator $`H,`$ where $`H`$ is an arbitrary, but fixed Schrödinger operator with a real potential $`qH_{loc}^1(^n)`$ such that $`\lambda _0(q)>0`$.
###### Lemma B.6.
The following estimates hold for any bounded open interval $`I`$ and any $`uH_{\mathrm{loc}}^1(I)`$:
(B.8)
$$u_{L^{\mathrm{}}(I)}\left|I\right|^{1/2}u_{L^2(I)}+u^{}_{L^1(I)},$$
and
(B.9)
$$u^{}_{L^1(I)}\left|I\right|^{1/2}Pu_{L^2(I)}+r_{L^2(I)}u_{L^2(I)}.$$
where $`|I|`$ means the length of the interval $`I`$.
###### Proof.
Both estimates are well-known but we provide the proofs for the convenience of the reader. We start with
$$u(x)u(y)=_y^xu^{}(s)𝑑s,x,yI,$$
which implies
$$|u(x)u(y)|u^{}_{L^1(I)},$$
hence
$$|u(x)||u(y)|+u^{}_{L^1(I)},$$
and
$$u_{L^{\mathrm{}}(I)}|u(y)|+u^{}_{L^1(I)},$$
for all $`yI`$. Integrating with respect to $`yI`$ and dividing by $`I`$, we get
$$u_{L^{\mathrm{}}(I)}|I|^1u_{L^1(I)}+u^{}_{L^1(I)}.$$
Applying the Cauchy-Schwarz inequality in the first term in the right hand side, we obtain (B.8).
From $`u^{}=Pu+ru`$, taking $`L^1`$-norms of both sides and using the Cauchy-Schwarz inequality we obtain
$$u^{}_{L^1(I)}Pu_{L^1(I)}+ru_{L^1(I)}|I|^{1/2}Pu_{L^2(I)}+r_{L^2(I)}u_{L^2(I)},$$
which proves (B.9). ∎
###### Lemma B.7.
Let us assume that $`qH_{loc}^1()`$ with $`\lambda _0(q)>0`$. Then
1) The operator $`\overline{P}T`$ is defined everywhere in $`L^2()`$ and $`\overline{P}T=\lambda _0(q)^{1/2}`$.
2) $`\overline{P}TP^{}=I`$ on $`𝔇(P^{})`$.
###### Proof.
1) Since $`P^{}(\overline{P}T)=(P^{}\overline{P})T=I`$, the operator $`\overline{P}T`$ is everywhere defined in $`L^2()`$. Also, for any $`uL^2()`$,
$$\overline{P}Tu^2=(\overline{P}Tu,\overline{P}Tu)=(TP^{}\overline{P}Tu,u)=(Tu,u)=T^{1/2}u^2,$$
so the first statement immediately follows. (In fact, the presentation $`\overline{P}T=UT^{1/2}`$, with $`U=\overline{P}T^{1/2}`$, is the polar decomposition of $`\overline{P}T`$.)
2) For any $`u,v𝔇(P^{}\overline{P})`$ we have
$$((\overline{P}TP^{})\overline{P}u,\overline{P}v)=(\overline{P}T(P^{}\overline{P})u,\overline{P}v)=(\overline{P}u,\overline{P}v),$$
so $`\overline{P}TP^{}u=u`$ for all $`u𝔇(P^{}\overline{P})`$. It remains to recall that $`𝔇(P^{}\overline{P})`$ is dense in $`L^2()`$. ∎
Heuristically, Green’s function $`G=G(x,y)`$ should be given by
(B.10)
$$G(x,y)=(T\delta (y))(x).$$
So it is expected to satisfy
(B.11)
$$(_x^2+q(x))G(x,y)=\delta (xy),$$
and be the Schwartz kernel of a bounded linear operator in $`L^2()`$. More precisely, we will prove:
###### Lemma B.8.
Let us assume that $`qH_{loc}^1()`$ with $`\lambda _0(q)>0`$. Then there exists a measurable, real-valued function $`G(x,y)`$ on $`\times `$ with the following properties:
(i) For any bounded interval $`I`$ in $``$, there is a positive constant $`C(I)`$ so that
$$\underset{xI}{sup}\left(G(x,y)^2𝑑y\right)^{1/2}C(I)$$
and the map $`xG(x,)`$ is Hölder continuous of order $`1/2`$ as a mapping from $`I`$ into $`L^2()`$.
(ii) For any $`fL^2()`$,
$$\left(Tf\right)(x)=G(x,y)f(y)𝑑y$$
(iii) $`G(x,y)=G(y,x)`$ almost everywhere.
(iv) For each fixed $`x`$, the function $`G(x,)`$ belongs to $`𝔇(\overline{P})`$.
(v) For any $`\phi 𝒞_0^{\mathrm{}}()`$ and any $`x,`$ $`(\overline{P}G(x,),\overline{P}\phi )=\phi (x).`$
###### Remark B.9.
Part (v) states that $`G(x,y)`$, viewed as a function of $`y`$ with $`x`$ as parameter, solves the equation $`L_qG(x,y)=\delta _x(y)`$ in distribution sense.
###### Proof of Lemma B.8.
In what follows, $`rL_{loc}^2()`$ is fixed and satisfies $`q=r^{}+r^2`$, $`I`$ denotes a bounded interval in $`,`$ and $`C(I)`$ denotes a generic constant depending on $`\left|I\right|`$ and $`r`$. Its value may vary from line to line.
We will make repeated use of the following observation, based on Lemma B.6. If $`\psi 𝒟(\overline{P})`$ and $`I`$ is a bounded interval, then by Lemma B.6,
(B.12)
$$\underset{xI}{sup}\left|\psi (x)\right|C(I)\left(\psi _{L^2(I)}+\overline{P}\psi _{L^2()}\right).$$
Using the boundedness of $`\psi `$ and the fact that $`\overline{P}\psi L^2()`$, we can then deduce that
(B.13) $`\psi ^{}_{L^2(I)}`$ $`\overline{P}\psi _{L^2()}+r_{L^2(I)}\psi _{L^{\mathrm{}}(I)}`$
$`C(I)\left(\psi _{L^2()}+\overline{P}\psi _{L^2()}\right).`$
Since, by Lemma B.7, $`\overline{P}T\psi _{L^2()}C\psi _{L^2()}`$, it follows from (B.12) and (B.13) that for any $`\psi L^2()`$, one has $`T\psi H_{\mathrm{loc}}^1()`$ with
(B.14)
$$\underset{xI}{sup}\left|(T\psi )(x)\right|C(I)\psi _{L^2()}$$
and
(B.15)
$$_I\left|\frac{d}{dx}(T\psi )(x)\right|^2𝑑xC(I)\psi _{L^2()}^2.$$
In particular, for each $`x`$, the map $`x\left(T\psi \right)(x)`$ is a bounded linear functional on $`L^2()`$. It follows from the Riesz representation theorem that there is an element $`G_x`$ of $`L^2()`$ with
$$(T\psi )(x)=(\psi ,G_x).$$
We claim that, also, the map $`IxG_xL^2()`$ is Hölder continuous of order $`1/2`$. To see this, we use (B.15) together with the Cauchy-Schwarz inequality to conclude that for $`x`$ and $`y`$ belonging to $`I`$,
$$\left|(T\psi )(x)(T\psi )(y)\right|C(I)\left|xy\right|^{1/2}\psi _{L^2()}$$
and thus
$$sup\left\{\right|(\psi ,G_xG_y)|:\psi L^2(),\psi _{L^2()}=1\}C(I)|xy|^{1/2}.$$
This proves the required Hölder continuity. It follows that the map $`xG_x`$ is a weakly measurable map from $`I`$ into $`L^2()`$ with $`G_x_{L^2()}`$ bounded uniformly in $`xI`$, so that $`xG_x`$ may be regarded as an element of the space $`L^2(I;L^2())`$ consisting of weakly measurable, square-integrable functions on $`I`$ taking values in $`L^2()`$. By Theorem III.11.17 of , there is a measurable function $`G_I(x,y)`$ on $`I\times `$ with the property that $`G_I(x,)=G_x`$ for every $`xI`$. As
$$(\phi ,T\psi )=_{I\times }\phi (x)G_I(x,y)\psi (y)𝑑y𝑑x$$
for any $`\phi L^{\mathrm{}}(I)`$ and $`\psi L^2()`$, it is easy to see that for any bounded intervals $`I`$ and $`J`$ with $`IJ`$, the restriction of $`G_J`$ to $`I\times `$ equals $`G_I`$ almost everywhere with respect to product measure on $`I\times `$. Taking a sequence of bounded intervals $`\left\{I_n\right\}`$ with $`I_n`$ as $`n\mathrm{}`$, we can construct a measurable function $`G`$ on $`\times `$ that obeys properties (i) and (ii). Property (iii) follows from the symmetry of $`T`$.
To prove property (iv), let $`\phi 𝒟(P^{})`$ and note that
$$(G_x,P^{}\phi )=(TP^{}\phi )(x).$$
By Lemma B.7, $`\overline{P}TP^{}\phi _{L^2()}\phi _{L^2()}`$ holds. Hence $`TP^{}\phi W_{\mathrm{loc}}^{1,1}()`$ and, for any bounded interval $`I`$,
$$\underset{xI}{sup}\left|\left(TP^{}\right)(x)\right|C(I)\phi _{L^2()}$$
so that
$$\left|(G_x,P^{}\phi )\right|C(I)\phi _{L^2()}$$
for any $`\phi 𝒟(P^{})`$. This shows that $`G_x𝒟(P^{})=𝒟(\overline{P})`$, proving (iv).
Finally, to prove (v), let $`\phi 𝒟(H)`$ and compute
$`\phi (x)`$ $`=(TH\phi )(x)`$
$`=(H\phi ,G_x)`$
$`=(\overline{P}\phi ,\overline{P}G_x).`$
Since $`𝒟(H)`$ is dense in $`𝒟(\overline{P})`$ and point evaluations are continuous in $`𝒟(\overline{P})`$, it follows that $`\phi (x)=(\overline{P}\phi ,\overline{P}G_x)`$ for all $`\phi 𝒟(\overline{P})`$. Since $`𝒞_0^{\mathrm{}}()𝒟(\overline{P})`$, (v) is proved. ∎
To construct positive solutions from Green’s function, we will need the following lemma.
###### Lemma B.10.
Suppose $`\lambda _0(q)>0`$ and that $`yH_{\mathrm{loc}}^1()`$, $`L_qy=0`$, and either
(i) $`yL^2(0,\mathrm{})`$ and $`PyL^2(0,\mathrm{})`$, or
(ii) $`yL^2(\mathrm{},0)`$ and $`PyL^2(\mathrm{},0)`$ .
Then, either $`y`$ has no zeros on $``$ or $`y`$ is identically zero on $``$.
###### Proof.
We will give the proof assuming (i) holds since the proof assuming (ii) holds is similar. Suppose that $`yH_{\mathrm{loc}}^1()`$ solves $`L_qy=0`$, (i) holds, and $`y(x_0)=0`$ for some $`x_0`$. We will assume without loss that $`x_0=0`$. By assumption (i), the function
$$w(x)=\{\begin{array}{cc}y(x),\hfill & 0x<\mathrm{}\hfill \\ 0,\hfill & x<0\hfill \end{array}$$
belongs to $`L^2()H_{\mathrm{loc}}^1()`$, and $`PwL^2()`$, hence $`wW_{\mathrm{loc}}^{1,1}()`$. It follows from Lemma B.1 that $`w𝔇(\overline{P}).`$
We claim that there is a sequence $`\left\{\phi _n\right\}`$ from $`𝒞_0^{\mathrm{}}()`$ with support contained in $`(0,\mathrm{})`$ so that $`\phi _nw`$ and $`P\phi _n\overline{P}w`$ in $`L^2()`$. If so then, on the one hand,
(B.16)
$$(P\phi _n,\overline{P}w)=0$$
since $`L_qy=0`$ and $`L_qy=L_qw`$ as distributions on $`(0,\mathrm{}).`$ Taking limits in (B.16) as $`n\mathrm{}`$ we have
(B.17)
$$(\overline{P}w,\overline{P}w)=0,$$
hence $`\overline{P}w=0`$. On the other hand, we have
(B.18)
$$(P\phi _n,P\phi _n)\lambda _0(q)\phi _n^2$$
where $`\lambda _0(q)>0`$. Taking limits as $`n\mathrm{}`$ in (B.18) and using (B.17), we conclude that $`w=0`$. It follows from the uniqueness of solutions to $`L_qy=0`$ with prescribed initial data that $`y=0`$ identically.
Thus, it remains to prove the existence of a sequence $`\left\{\phi _n\right\}`$ with the claimed properties. First, we show that $`w`$ may be approximated by functions which vanish identically near $`x=0`$. Let $`\chi 𝒞^{\mathrm{}}()`$ with $`0\chi (x)1`$, $`\chi (x)=0`$ for $`x1`$, $`\chi (x)=1`$ for $`x2`$, and $`|\chi ^{}(x)|2`$ for all $`x`$. Let $`\chi _\epsilon (x)=\chi (x/\epsilon )`$. The functions $`w_\epsilon (x)=\chi _\epsilon (x)w(x)`$ converge to $`w`$ in $`L^2()`$ as $`\epsilon 0`$ by dominated convergence. We claim that, also, $`Pw_\epsilon \overline{P}w`$ in $`L^2()`$. To see this, compute
(B.19)
$$Pw_\epsilon =\chi _\epsilon ^{}(x)w(x)+\chi _\epsilon (x)(\overline{P}w)(x).$$
The second term in the right-hand side of (B.19) converges to $`\overline{P}w`$ in $`L^2()`$ by dominated convergence while the first one converges to $`0`$ in $`L^2()`$ by the following reasons. Observing that
$$\chi _\epsilon ^{}(x)^2\left|w(x)\right|^2𝑑x\frac{4}{\epsilon ^2}_\epsilon ^{2\epsilon }\left(_0^x\left|w^{}(t)\right|𝑑t\right)^2𝑑x4_0^{2\epsilon }\left|w^{}(t)\right|^2𝑑t,$$
we conclude
$$\chi _\epsilon ^{}(x)^2\left|w(x)\right|^2𝑑x0\text{ as }\epsilon 0.$$
Thus $`Pw_\epsilon \overline{P}w`$ in $`L^2()`$.
Letting $`\epsilon =1/n`$, the function $`w_{1/n}`$ has support in $`[1/n,\mathrm{})`$. We can use smooth cut-off functions and Friedrichs mollifiers as in the proof of Lemma B.1 to find a $`𝒞_0^{\mathrm{}}`$ function $`\phi _n`$ with support in $`[1/(2n),\mathrm{})`$ so that $`\phi _nw_{1/n}<1/n`$ and $`P\phi _n\overline{P}w_{1/n}<1/n`$. In this way we obtain a sequence $`\left\{\phi _n\right\}`$ from $`𝒞_0^{\mathrm{}}(0,\mathrm{})`$ so that $`\phi _nw0`$ and $`P\phi _n\overline{P}w0`$ as $`n\mathrm{}`$. ∎
As an application of the results obtained in this appendix, we give an alternative proof of Proposition 5.8
Proof of Proposition 5.8. We claim that there exists an $`x`$ so that $`y`$ $`G(x,y)`$ does not vanish identically on $`(x,\mathrm{})`$. If not then $`G(x,y)=0`$ for all $`(x,y)`$ with $`y>x`$ and hence, by Lemma B.8(iii), for all $`yx`$. Therefore $`G(x,y)=0`$ a.e., a contradiction. Now choose such an $`x`$. Then the function $`\psi _+(y)=G(x,y)`$ for $`y>x`$ is not identically zero on $`(x,\mathrm{})`$. From Lemma B.8(i), (iv), and (v), $`\psi _+(y)L^2(x,\mathrm{})`$, $`P\psi _+L^2(x,\mathrm{})`$, and $`L_q\psi _+=0`$ for $`y>x`$. Let $`Q`$ be an antiderivative of $`q`$ and let $`\{y_+,u_+\}`$ be the unique solution to the system (2.7) with initial data $`y_+(x+1)=\psi _+(x+1)`$ and $`\left(u_+\right)(x+1)=\left(\psi _+Q\psi _+\right)(x+1)`$. Then $`y_+`$ coincides with $`\psi _+`$ on $`(x,\mathrm{})`$, so $`y_+`$ and $`Py_+`$ belong to $`L^2(0,\mathrm{})`$. It follows from Lemma B.10 that $`y_+`$ has no zeros, so by changing signs if necessary we conclude that $`y_+L^2(0,\mathrm{})`$ and $`y_+`$ is strictly positive on $``$. A similar construction considering the function $`\psi _{}(y)=G(x,y)`$ for some $`x`$ and $`y<x`$ leads to a strictly positive solution $`y_{}`$ of $`L_qy=0`$ with $`y_{}L^2(\mathrm{},0)`$. If $`y_+`$ and $`y_{}`$ were linearly dependent, then after multiplying one of them by an appropriate constant, we would obtain a function $`\psi `$ in the domain of $`H`$, $`\psi `$ not identically zero, with $`H\psi =0`$, which is impossible since $`\lambda _0(q)>0`$. Thus $`y_+`$ and $`y_{}`$ are linearly independent. $`\mathrm{}`$
## Appendix C Related Work
The Miura map was introduced by Miura , and played an important role in the search of integrals of motion for the Korteweg-de Vries equation. Miura discovered that his map takes smooth solutions of mKdV to smooth solutions of KdV. Hence it can serve as a tool to derive results on the initial value problem for KdV from results on the initial value problem for mKdV – see e.g. . Despite the fact that the Miura map is not one-to-one, when considered, for example, as a map between appropriate Sobolev spaces, it is also possible to use it to derive results for the initial value problem of mKdV from results of the initial value problem of KdV – see e.g. , , .
Miura map on the circle: Motivated by earlier work of Ambrosetti and Prodi on certain nonlinear elliptic boundary value problems, McKean and Scovel studied - among other nonlinear maps - the Miura map on the circle $`𝕋.`$ They exhibited a global fold structure for the Miura map when viewed as a map from $`H^1(𝕋)`$ to $`L^2(𝕋).`$ For further results in this direction, see Bueno and Tomei . Later, Korotyaev , and Kappeler and Topalov extended the global fold picture to the Miura map from periodic functions in $`L^2(𝕋)`$ to $`H^1(𝕋)`$. Kappeler and Topalov proved existence and well-posedness of solutions to the mKdV equation with initial data in $`L^2(𝕋)`$ , using and their results on the initial value problem for the periodic KdV equation established in .
Miura map on the line: On the line, the Miura map and related topics have also been investigated extensively, and not exclusively with a view towards applications for solving the initial value problem of KdV or mKdV.
$``$ Positive solutions of Schrödinger equations or more generally of second order elliptic equations – in particular in connection with spectral properties of the corresponding operators – have been extensively studied in various settings. We only mention the result, referred to as Allegretto-Piepenbrink theorem in , Theorem 2.12 or in , section C.8. This theorem states that for potentials $`qL_{\mathrm{loc}}^1(^n),`$ satisfying some additional conditions, $`(\mathrm{\Delta }+q)u=\lambda u`$ has a nonzero solution $`u`$ (in the sense that $`uW_{\mathrm{loc}}^{2,1}(^n)`$ and $`quL_{\mathrm{loc}}^1(^n)`$), which is nonnegative everywhere, if and only if $`inf(spec`$($`\mathrm{\Delta }+q))\lambda `$. See or for further details and references to the papers of Allegretto and of Piepenbrink as well as additional references. In the one-dimensional case at hand, the equivalence of the statements (ii) and (iii) of Theorem 1.1 for potentials $`qL_{\mathrm{loc}}^1()`$ is well known – see , Theorems XI.6.1 and XI.6.2 and Corollary XI.6.1, , Appendix 1, or , Theorem 3.1. Thus Theorem 1.1 as stated above shows in particular that this equivalence continues to hold for $`qH_{\mathrm{loc}}^1()`$.
$``$ With regard to the characterization of the image of the Miura map $`B_0`$ , we mention the result of Ablowitz et. al. which characterizes the image of Schwartz space by $`B_0`$ in terms of the scattering data of these potentials as well as the result of Tsutsumi , stating that any finite, positive Radon measure is in the image of the Miura map $`B_0:L^2()H^1().`$ Further, the case where $`q`$ is continuous is treated by Hartman , Chapter XI.7, Lemma 7.1. The result stated in Theorem 1.2 sharpens all these results and puts them into a broader perspective.
$``$ The dichotomy described in Theorem 1.3 has another interpretation which does not involve the Miura map at all: Murata , Appendix 1, describes the dichotomy stated in Theorem 1.3 – again for potentials in $`qL_{\mathrm{loc}}^1()`$ – in terms of the notion of subcritical, critical, and supercritical potentials, where in his terminology (i) $`q`$ is called subcritical if $`L_q`$ has a positive Green’s function, (ii) $`q`$ is called critical if $`L_q0`$ and does not have a positive Green’s function , and (iii) $`q`$ is called supercritical if $`L_q`$ is not nonnegative. See Simon for an alternative notion of subcritical and critical potentials. In , Theorem A.5, Murata shows that (i) $`qL_{\mathrm{loc}}^1()`$ is subcritical iff $`L_qy=0`$ admits two linearly independent positive solutions in $`W_{\mathrm{loc}}^{2,1}()`$ and that (ii) $`qL_{\mathrm{loc}}^1()`$ is critical iff $`L_qy=0`$ has up to scaling one positive solution $`yW_{\mathrm{loc}}^{2,1}().`$ These results of Murata for one-dimensional Schrödinger operators were later extended by Gesztesy and Zhao , Theorem 3.6, to more general Sturm-Liouville operators. Our results on the dichotomy for Schrödinger operators obtained in this paper extend the results of Murata (and of Gesztesy and Zhao) in two directions. First, we consider potentials which are real-valued distributions in a Sobolev space with negative index of smoothness $`\beta 1`$. Hence they are not necessarily functions. Second, we describe geometric aspects of the dichotomy: see Theorem 1.3 and Theorem 5.6.
Schrödinger operators with singular potentials: Recently, the operators $`L_q`$, considered on an interval $`(a,b)`$, with potential $`q`$ in a Sobolev space with negative index of smoothness, have been studied by various authors. In particular, we mention the paper where different approaches to define the operator $`L_q`$ are discussed in detail and asymptotics for the eigenvalues and eigenfunctions of these operators are obtained. See also , , , , , , , as well as for further references.
Initial value problem for KdV: The initial value problem for KdV on the line has been extensively studied. We only mention that, based on the works of Bourgain , , it has been proved by Kenig-Ponce-Vega that KdV is locally uniformly $`C^0`$ well-posed on $`H^s()`$ for $`s>3/4`$ and later, by Colliander, Keel, Staffilani, Takaoka and Tao , that KdV is globally uniformly $`C^0`$ well-posed on $`H^s()`$ for $`s>3/4`$. An existence result for the limiting case $`s=3/4`$ has been obtained by Christ, Colliander and Tao . Beside the work of Tsutsumi already mentioned above on solutions of KdV with positive Radon measures as initial data, it has been shown in that for measures of bounded variation with sufficient decay at infinity as initial data, there exists a *classical* solution for $`t>0`$.
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# A generating function method for the average-case analysis of DPLL
## 1 Introduction and main results.
Many efforts have been devoted to the study of the performances of the Davis-Putnam-Loveland-Logemann (DPLL) procedure , and more generally, resolution proof complexity for combinatorial problems with randomly generated instances. Two examples are random $`k`$-Satisfiability ($`k`$-SAT), where an instance $``$ is a uniformly and randomly chosen set of $`M=\alpha N`$ disjunctions of $`k`$ literals built from $`N`$ Boolean variables and their negations (with no repetition and no complementary literals), and random graph $`k`$-Coloring ($`k`$-COL), where an instance $``$ is an Erdős-Rényi random graph from $`G(N,p=c/N)`$ i.e. with average vertex degree $`c`$.
Originally, efforts were concentrated on the random width distribution for $`k`$-SAT, where each literal appear with a fixed probability. Franco, Purdom and collaborators showed that simplified versions of DPLL had polynomial average-case complexity in this case, see for reviews. It was then recognized that the fixed clause length ensemble might provide harder instances for DPLL . Chvátal and Szemerédi indeed showed that DPLL proof size is w.h.p. exponentially large (in $`N`$ at fixed ratio $`\alpha `$) for an unsatisfiable instance . Later on, Beame et al. showed that the proof size was w.h.p. bounded from above by $`2^{cN/\alpha }`$ (for some constant $`c`$), a decreasing function of $`\alpha `$. As for the satisfiable case, Frieze and Suen showed that backtracking is irrelevant at small enough ratios $`\alpha `$ ($`3.003`$ with the Generalized Unit Clause heuristic, to be defined below) , allowing DPLL to find satisfying assignment in polynomial (linear) time. Achlioptas, Beame and Molloy proved that, conversely, at ratios smaller than the generally accepted satisfiability threshold, DPLL takes w.h.p. exponential time to find a satisfying assignment . Altogether these results provide explanations for the ‘easy-hard-easy’ (or, more precisely, ‘easy-hard-less hard’) pattern of complexity experimentally observed when running DPLL on random 3-SAT instances .
A precise calculation of the average size of the search space explored by DPLL (and #DPLL, a version of the procedure solving the enumeration problems #SAT and #COL) as a function of the parameters $`N`$ and $`\alpha `$ or $`c`$ is difficult due to the statistical correlations between branches in the search tree resulting from backtracking. Heuristic derivations were nevertheless proposed by Cocco and Monasson based on a ‘dynamic annealing’ assumption . Hereafter, using the linearity of expectation, we show that ‘dynamic annealing’ turns not to be an assumption at all when the expected tree size is concerned.
We first illustrate the approach, based on the use of recurrence relations for the generating functions of the number of nodes at a given height in the tree, on the random $`k`$-SAT problem and the simple Unit Clause (UC) branching heuristic where unset variables are chosen uniformly at random and assigned to True or False uniformly at random . Consider the following counting algorithm Procedure #DPLL-UC\[$``$,A,S\]
Call $`_A`$ what is left from instance $``$ given partial variable assignment $`A`$;
1. If $`_A`$ is empty, $`SS+2^{N|A|}`$, Return; (Solution Leaf)
2. If there is an empty clause in $`_A`$, Return; (Contradiction Leaf)
3. If there is no empty clause in $`_A`$, let $`\mathrm{\Gamma }_1=\{`$1-clauses$`_A\}`$,
if $`\mathrm{\Gamma }_1\mathrm{}`$, pick any 1-clause, say, $`\mathrm{}`$, and call DPLL\[$``$,A$`\mathrm{}`$\]; (unit-propagation)
if $`\mathrm{\Gamma }_1=\mathrm{}`$, pick up an unset literal uniformly at random, say, $`\mathrm{}`$, and call
DPLL\[$``$,A$`\mathrm{}`$\], then DPLL\[$``$,A$`\overline{\mathrm{}}`$\] ; (variable splitting)
End;
#DPLL-UC, called with $`A=\mathrm{}`$ and $`S=0`$, returns the number $`S`$ of solutions of the instance $``$; the history of the search can be summarized as a search tree with leaves marked with solution or contradiction labels. As the instance to be treated and the sequence of operations done by #DPLL-UC are stochastic, so are the numbers $`L_S`$ and $`L_C`$ of solution and contradiction leaves respectively.
###### Theorem 1.1
Let $`k3`$ and $`\mathrm{\Omega }(t,\alpha ,k)=t+\alpha \mathrm{log}_2\left(1{\displaystyle \frac{k}{2^k}}t^{k1}+{\displaystyle \frac{k1}{2^k}}t^k\right)`$. The expectations of the numbers of solution and contradiction leaves in the #DPLL-UC search tree of random $`k`$-SAT instances with $`N`$ variables and $`\alpha N`$ clauses are, respectively, $`L_S(N,\alpha ,k)=2^{N\omega _S(\alpha ,k)+o(N)}`$ with $`\omega _S(\alpha ,k)=\mathrm{\Omega }(1,\alpha ,k)`$ and $`L_C(N,\alpha ,k)=2^{N\omega _C(\alpha ,k)+o(N)}with\omega _C(\alpha ,k)=\underset{t[0;1]}{\mathrm{max}}\mathrm{\Omega }(t,\alpha ,k)`$.
An immediate consequence of Theorem 1 is that the expectation value of the total number of leaves, $`L_S+L_C`$, is $`2^{N\omega _C(\alpha ,k)+o(N)}`$. This result was first found by Méjean, Morel and Reynaud in the particular case $`k=3`$ and for ratios $`\alpha >1`$ . Our approach not only provides a much shorter proof, but can also be easily extended to other problems and more sophisticated heuristics, see Theorems 2 and 3 below. In addition, Theorem 1 provides us with some information about the expected search tree size of the decision procedure DPLL-UC, corresponding to #DPLL-UC with Line 1 replaced with: If $`_A`$ is empty, output Satisfiable; Halt.
###### Corollary 1
Let $`\alpha >\alpha _u(k)`$, the root of $`\omega _C(\alpha ,k)=2+\alpha \mathrm{log}_2(12^k)`$ e.g. $`\alpha _u(3)=10.1286\mathrm{}`$. The average size of DPLL-UC search trees for random $`k`$-SAT instances with $`N`$ variables and $`\alpha N`$ clauses equals $`2^{N\omega _C(\alpha ,k)+o(N)}`$.
Functions $`\omega _S,\omega _C`$ are shown in Figure 1 in the $`k=3`$ case. They coincide and are equal to $`1\alpha \mathrm{log}_2(8/7)`$ for $`\alpha <\alpha ^{}=4.56429\mathrm{}`$, while $`\omega _C>\omega _S`$ for $`\alpha >\alpha ^{}`$. In other words, for $`\alpha >\alpha ^{}`$, most leaves in #DPLL-UC trees are contradiction leaves, while for $`\alpha <\alpha ^{}`$, both contradiction and solution leaf numbers are (to exponential order in $`N`$) of the same order. As for DPLL-UC trees, notice that $`\omega _C(\alpha ,k){\displaystyle \frac{2\mathrm{ln}2}{3\alpha }}={\displaystyle \frac{0.46209\mathrm{}}{\alpha }}`$. This behaviour agrees with Beame et al.’s result ($`\mathrm{\Theta }(1/\alpha )`$) for the average resolution proof complexity of unsatisfiable instances . Corollary 1 shows that the expected DPLL tree size can be estimated for a whole range of $`\alpha `$; we conjecture that the above expression holds for ratios smaller than $`\alpha _u`$ i.e. down to $`\alpha ^{}`$ roughly. For generic $`k3`$, we have $`\omega _C(\alpha ,k){\displaystyle \frac{k2}{k1}}\left({\displaystyle \frac{2^k\mathrm{ln}2}{k(k1)\alpha }}\right)^{1/(k2)}`$; the decrease of $`\omega _C`$ with $`\alpha `$ is therefore slower and slower as $`k`$ increases.
So far, no expression for $`\omega `$ has been obtained for more sophisticated heuristics than UC. We consider the Generalized Unit Clause (GUC) heuristic where the shortest clauses are preferentially satisfied. The associated decision procedure, DPLL-GUC, corresponds to DPLL-UC with Line 3 replaced with: Pick a clause uniformly at random among the shortest clauses, and a literal, say, $`\mathrm{}`$, in the clause; call DPLL\[$``$,A$`\mathrm{}`$\], then DPLL\[$``$,A$`\overline{\mathrm{}}`$\].
###### Theorem 1.2
Define $`m(x_2)=\frac{1}{2}(1+\sqrt{1+4x_2})2x_2`$, $`y_3(y_2)`$ the solution of the ordinary differential equation $`dy_3/dy_2=3(1+y_22y_3)/(2m(y_2))`$ such that $`y_3(1)=1`$, and
$$\omega ^g(\alpha )=\underset{\frac{3}{4}<y_21}{\mathrm{max}}\left[_{y_2}^1\frac{dz}{m(z)}\mathrm{log}_2\left(2z+m(z)\right)\mathrm{exp}\left(_z^1\frac{dw}{m(w)}\right)+\alpha \mathrm{log}_2y_3(y_2)\right].$$
Let $`\alpha >\alpha _u^g=10.2183\mathrm{}`$, the root of $`\omega ^g(\alpha )+\alpha \mathrm{log}_2(8/7)=2`$. The expected size of DPLL-GUC search tree for random 3-SAT instances with $`N`$ variables and $`\alpha N`$ clauses is $`2^{N\omega ^g(\alpha )+o(N)}`$.
Notice that, at large $`\alpha `$, $`\omega ^g(\alpha ){\displaystyle \frac{3+\sqrt{5}}{6\mathrm{ln}2}}\left[\mathrm{ln}\left({\displaystyle \frac{1+\sqrt{5}}{2}}\right)\right]^2{\displaystyle \frac{1}{\alpha }}={\displaystyle \frac{0.29154\mathrm{}}{\alpha }}`$ in agreement with the $`1/\alpha `$ scaling established in . Furthermore, the multiplicative factor is smaller than the one for UC, showing that DPLL-GUC is more efficient than DPLL-UC in proving unsatisfiability. A third application is the analysis of the counterpart of GUC for the random 3-COL problem. The version of DPLL we have analyzed operates as follows . Initially, each vertex is assigned a list of 3 available colors. In the course of the procedure, a vertex, say, $`v`$, with the smallest number of available colors, say, $`j`$, is chosen at random and uniformly. DPLL-GUC then removes $`v`$, and successively branches to the $`j`$ color assignments corresponding to removal of one of the $`j`$ colors of $`v`$ from the lists of the neighbors of $`v`$. The procedure backtracks when a vertex with no color left is created (contradiction), or no vertex is left (a proper coloring is found).
###### Theorem 1.3
Define $`\omega ^h(c)=\underset{0<t<1}{\mathrm{max}}\left[{\displaystyle \frac{c}{6}}t^2{\displaystyle \frac{c}{3}}t(1t)\mathrm{ln}2+\mathrm{ln}\left(3e^{2ct/3}\right)\right]`$.
Let $`c>c_u^h=13.1538\mathrm{}`$, the root of $`\omega ^h(c)+\frac{c}{6}=2\mathrm{ln}3`$. The expected size of DPLL-GUC search tree for deciding 3-COL on random graphs from $`G(N,c/N)`$ with $`N`$ vertices is $`e^{N\omega ^h(c)+o(N)}`$.
Asymptotically, $`\omega ^h(c){\displaystyle \frac{3\mathrm{ln}2}{2c^2}}={\displaystyle \frac{1.0397\mathrm{}}{c^2}}`$ in agreement with Beame et al.’s scaling ($`\mathrm{\Theta }(1/c^2)`$) . An extension of Theorem 3 to higher values of the number $`k`$ of colors gives $`\omega ^h(c,k){\displaystyle \frac{k(k2)}{k1}}\left[{\displaystyle \frac{2\mathrm{ln}2}{k1}}\right]^{1/(k2)}c^{(k1)/(k2)}`$. This result is compatible with the bounds derived in , and suggests that the $`\mathrm{\Theta }(c^{(k1)/(k2)})`$ dependence could hold w.h.p. (and not only in expectation).
## 2 Recurrence equation for #DPLL-UC search tree
Let $``$ be an instance of the 3-SAT problem defined over a set of $`N`$ Boolean variables $`X`$. A partial assignment $`A`$ of length $`T(N)`$ is the specification of the truth values of $`T`$ variables in $`X`$. We denote by $`_A`$ the residual instance given $`A`$. A clause $`c_A`$ is said to be a $`\mathrm{}`$-clause with $`\mathrm{}\{0,1,2,3\}`$ if the number of false literals in $`c`$ is equal to $`3\mathrm{}`$. We denote by $`C_{\mathrm{}}(_A)`$ the number of $`\mathrm{}`$-clauses in $`_A`$. The instance $``$ is said to be satisfied under $`A`$ if $`C_{\mathrm{}}(_A)=0`$ for $`\mathrm{}=0,1,2,3`$, unsatisfied (or violated) under $`A`$ if $`C_0(_A)1`$, undetermined under $`A`$ otherwise. The clause vector of an undetermined or satisfied residual instance $`_A`$ is the three-dimensional vector $`\stackrel{}{C}`$ with components $`C_1(_A),C_2(_A),C_3(_A)`$. The search tree associated to an instance $``$ and a run of #DPLL is the tree whose nodes carry the residual assignments $`A`$ considered in the course of the search. The height $`T`$ of a node is the length of the attached assignment.
It was shown by Chao and Franco that, during the first descent in the search tree i.e. prior to any backtracking, the distribution of residual instances remains uniformly random conditioned on the numbers of $`\mathrm{}`$-clauses. This statement remains correct for heuristics more sophisticated than UC e.g. GUC, SC<sub>1</sub> , and was recently extended to splitting heuristics based on variable occurrences by Kaporis, Kirousis and Lalas . Clearly, in this context, uniformity is lost after backtracking enters into play (with the exception of Suen and Frieze’s analysis of a limited version of backtracking ). Though this limitation appears to forbid (and has forbidden so far) the extension of average-case studies of backtrack-free DPLL to full DPLL with backtracking, we point out here that it is not as severe as it looks. Indeed, let us forget about how #DPLL or DPLL search tree is built and consider its final state. We refer to a branch (of the search tree) as the shortest path from the root node (empty assignment) to a leaf. The two key remarks underlying the present work can be informally stated as follows. First, the expected size of a #DPLL search tree can be calculated from the knolwedge of the statistical distribution of (residual instances on) a single branch; no characterization of the correlations between distinct branches in the tree is necessary. Secondly, the statistical distribution of (residual instances on) a single branch is simple since, along a branch, uniformity is preserved (as in the absence of backtracking). More precisely,
###### Lemma 1 (from Chao & Franco )
Let $`_A`$ be a residual instance attached to a node $`A`$ at height $`T`$ in a #DPLL-UC search tree produced from an instance $``$ drawn from the random 3-SAT distribution. Then the set of $`\mathrm{}`$-clauses in $`_A`$ is uniformly random conditioned on its size $`C_{\mathrm{}}(_A)`$ and the number $`NT`$ of unassigned variables for each $`\mathrm{}\{0,1,2,3\}`$.
###### Proof
the above Lemma is an immediate application of Lemma 3 in Achlioptas’ Card Game framework which establishes uniformity for algorithms (a) ‘pointing to a particular card (clause)’, or (b) ’naming a variable that has not yet been assigned a value’ (Section 2.1 in Ref. ). The operation of #DPLL-UC along a branch precisely amounts to these two operations: unit-propagation relies on action (a), and variable splitting on (b). ∎
Lemma 1 does not address the question of uniformity among different branches. Residual instances attached to two (or more) nodes on distinct branches in the search tree are correlated. However, these correlations can be safely ignored in calculating the average number of residual instances, in much the same way as the average value of the sum of correlated random variables is simply the sum of their average values.
###### Proposition 1
Let $`L(\stackrel{}{C},T)`$ be the expectation of the number of undetermined residual instances with clause vector $`\stackrel{}{C}`$ at height $`T`$ in #DPLL-UC search tree, and $`G(x_1,x_2,x_3;T)={\displaystyle \underset{\stackrel{}{C}}{}}x_1^{C_1}x_2^{C_2}x_3^{C_3}L(\stackrel{}{C},T)`$ its generating function. Then, for $`0T<N`$,
$`G(x_1,x_2,x_3;T+1)`$ $`=`$ $`{\displaystyle \frac{1}{f_1}}G(f_1,f_2,f_3;T)+\left(2{\displaystyle \frac{1}{f_1}}\right)G(0,f_2,f_3;T)`$ (1)
$``$ $`2G(0,0,0;T)`$
where $`f_1,f_2,f_3`$ stand for the functions $`f_1^{(T)}(x_1)=x_1+\frac{1}{2}\mu (12x_1)`$, $`f_2^{(T)}(x_1,x_2)=x_2+\mu (x_1+12x_2)`$, $`f_3^{(T)}(x_2,x_3)=x_3+\frac{3}{2}\mu (x_2+12x_3)`$, and $`\mu =1/(NT)`$. The generating function $`G`$ is entirely defined from recurrence relation (1) and the initial condition $`G(x_1,x_2,x_3;0)=\left(x_3\right)^{\alpha N}`$.
###### Proof
Let $`\delta _n`$ denote the Kronecker function ($`\delta _n=1`$ if $`n=0`$, $`\delta _n=0`$ otherwise), $`B_n^{m,q}=\left(\genfrac{}{}{0pt}{}{m}{n}\right)q^n(1q)^{mn}`$ the binomial distribution. Let $`A`$ be a node at height $`T`$, and $`_A`$ the attached residual instance. Call $`\stackrel{}{C}`$ the clause vector of $`_A`$. Assume first that $`C_11`$. Pick up one 1-clause, say, $`\mathrm{}`$. Call $`z_j`$ the number of $`j`$-clauses that contain $`\overline{\mathrm{}}`$ or $`\mathrm{}`$ (for $`j=1,2,3`$). From Lemma 1, the $`z_j`$’s are binomial variables with parameter $`j/(NT)`$ among $`C_j\delta _{j1}`$ (the 1-clause that is satisfied through unit-propagation is removed). Among the $`z_j`$ clauses, $`w_{j1}`$ contained $`\overline{\mathrm{}}`$ and are reduced to $`(j1)`$-clauses, while the remaining $`z_jw_{j1}`$ contained $`\mathrm{}`$ and are satisfied and removed. From Lemma 1 again, $`w_{j1}`$ is a binomial variable with parameter $`1/2`$ among $`z_j`$. The probability that the instance produced has no empty clause ($`w_0=0`$) is $`B_0^{z_1,\frac{1}{2}}=2^{z_1}`$. Thus, setting $`\mu =\frac{1}{NT}`$,
$`M_P`$ $`[\stackrel{}{C}^{},\stackrel{}{C};T]={\displaystyle \underset{z_3=0}{\overset{C_3}{}}}B_{z_3}^{C_3,3\mu }{\displaystyle \underset{w_2=0}{\overset{z_3}{}}}B_{w_2}^{z_3,\frac{1}{2}}{\displaystyle \underset{z_2=0}{\overset{C_2}{}}}B_{z_2}^{C_2,2\mu }{\displaystyle \underset{w_1=0}{\overset{z_2}{}}}B_{w_1}^{z_2,\frac{1}{2}}`$
$`\times {\displaystyle \underset{z_1=0}{\overset{C_11}{}}}B_{z_1}^{C_11,\mu }{\displaystyle \frac{1}{2^{z_1}}}\delta _{C_3^{}(C_3z_3)}\delta _{C_2^{}(C_2z_2+w_2)}\delta _{C_1^{}(C_11z_1+w_1)}`$
expresses the probability that a residual instance at height $`T`$ with clause vector $`\stackrel{}{C}`$ gives rise to a (non-violated) residual instance with clause vector $`\stackrel{}{C}^{}`$ at height $`T+1`$ through unit-propagation. Assume now $`C_1=0`$. Then, a yet unset variable is chosen and set to True or False uniformly at random. The calculation of the new vector $`\stackrel{}{C}^{}`$ is identical to the unit-propagation case above, except that: $`z_1=w_0=0`$ (absence of 1-clauses), and two nodes are produced (instead of one). Hence,
$`M_{UC}[\stackrel{}{C}^{},\stackrel{}{C};T]`$ $`=`$ $`2{\displaystyle \underset{z_3=0}{\overset{C_3}{}}}B_{z_3}^{C_3,3\mu }{\displaystyle \underset{w_2=0}{\overset{z_3}{}}}B_{w_2}^{z_3,\frac{1}{2}}{\displaystyle \underset{z_2=0}{\overset{C_2}{}}}B_{z_2}^{C_2,2\mu }{\displaystyle \underset{w_1=0}{\overset{z_2}{}}}B_{w_1}^{z_2,\frac{1}{2}}`$
$`\times \delta _{C_3^{}(C_3z_3)}\delta _{C_2^{}(C_2z_2+w_2)}\delta _{C_1^{}w_1}`$
expresses the expected number of residual instances at height $`T+1`$ and with clause vector $`\stackrel{}{C}^{}`$ produced from a residual instance at height $`T`$ and with clause vector $`\stackrel{}{C}`$ through UC branching.
Now, consider all the nodes $`A_i`$ at height $`T`$, with $`i=1,\mathrm{},`$. Let $`o_i`$ be the operation done by #DPLL-UC on $`A_i`$. $`o_i`$ represents either unit-propagation (literal $`\mathrm{}_i`$ set to True) or variable splitting (literals $`\mathrm{}_i`$ set to T and F on the descendent nodes respectively). Denoting by $`𝐄_Y(X)`$ the expectation value of a quantity $`X`$ over variable $`Y`$, $`L(\stackrel{}{C}^{};T+1)=𝐄_{,\{A_i,o_i\}}\left({\displaystyle \underset{i=1}{\overset{}{}}}[\stackrel{}{C}^{};A_i,o_i]\right)`$ where $``$ is the number (0, 1 or 2) of residual instances with clause vector $`\stackrel{}{C}^{}`$ produced from $`A_i`$ after #DPLL-UC has carried out operation $`o_i`$. Using the linearity of expectation, $`L(\stackrel{}{C}^{};T+1)=𝐄_{}\left({\displaystyle \underset{i=1}{\overset{}{}}}𝐄_{\{A_i,o_i\}}\left([\stackrel{}{C}^{};A_i,o_i]\right)\right)=𝐄_{}\left({\displaystyle \underset{i=1}{\overset{}{}}}M[\stackrel{}{C}^{},\stackrel{}{C}_i;T]\right)`$
where $`\stackrel{}{C}_i`$ is the clause vector of the residual instance attached to $`A_i`$, and $`M[\stackrel{}{C}^{},\stackrel{}{C};T]=\left(1\delta _{C_1}\right)M_P[\stackrel{}{C}^{},\stackrel{}{C};T]+\delta _{C_1}M_{UC}[\stackrel{}{C}^{},\stackrel{}{C};T]`$. Gathering assignments with identical clause vectors gives the reccurence relation $`L(\stackrel{}{C}^{},T+1)=`$ $`{\displaystyle \underset{\stackrel{}{C}}{}}M[\stackrel{}{C}^{},\stackrel{}{C};T]L(\stackrel{}{C},T)`$. Recurrence relation (1) for the generating function is an immediate consequence. The initial condition over $`G`$ stems from the fact that the instance is originally drawn from the random 3-SAT distribution, $`L(\stackrel{}{C};0)=\delta _{C_1}\delta _{C_2}\delta _{C_3\alpha N}`$.∎
## 3 Asymptotic analysis and application to DPLL-UC
The asymptotic analysis of $`G`$ relies on the following technical lemma:
###### Lemma 2
Let $`\gamma (x_2,x_3,t)=(1t)^3x_3+\frac{3t}{2}(1t)^2x_2+\frac{t}{8}(123t2t^2)`$, with $`t]0;1[`$ and $`x_2,x_3>0`$. Define $`S_0(T){\displaystyle \underset{H=0}{\overset{T}{}}}2^{TH}G(0,0,0;H)`$. Then, in the large $`N`$ limit, $`S_0([tN])2^{N(t+\alpha \mathrm{log}_2\gamma (0,0,t))+o(N)}`$ and $`G(\frac{1}{2},x_2,x_3;[tN])=2^{N(t+\alpha \mathrm{log}_2\gamma (x_2,x_3,t))+o(N)}`$.
Due to space limitations, we give here only some elements of the proof. The first step in the proof is inspired by Knuth’s kernel method : when $`x_1=\frac{1}{2}`$, $`f_1=\frac{1}{2}`$ and recurrence relation (1) simplifies and is easier to handle. Iterating this equation then allows us to relate the value of $`G`$ at height $`T`$ and coordinates $`(\frac{1}{2},x_2,x_3)`$ to the (known) value of $`G`$ at height 0 and coordinates $`(\frac{1}{2},y_2,y_3)`$ which are functions of $`x_2,x_3,T,N`$, and $`\alpha `$. The function $`\gamma `$ is the value of $`y_3`$ when $`T,N`$ are sent to infinity at fixed ratio $`t`$. The asymptotic statement about $`S_0(T)`$ comes from the previous result and the fact that the dominant terms in the sum defining $`S_0`$ are the ones with $`H`$ close to $`T`$.
###### Proposition 2
Let $`L_C(N,T,\alpha )`$ be the expected number of contradiction leaves of height $`T`$ in the #DPLL-UC resolution tree of random 3-SAT instances with $`N`$ variables and $`\alpha N`$ clauses, and $`ϵ>0`$. Then, for $`t[ϵ;1ϵ]`$ and $`\alpha >0`$, $`\mathrm{\Omega }(t,\alpha ,3){\displaystyle \frac{1}{N}}\mathrm{log}_2L_C(N,[tN],\alpha )+o(1)\underset{h[ϵ,;t]}{\mathrm{max}}\mathrm{\Omega }(h,\alpha ,3)`$ where $`\mathrm{\Omega }`$ is defined in Theorem 1.
Observe that a contradiction may appear with a positive (and non–exponentially small in $`N`$) probability as soon as two 1-clauses are present. These 1-clauses will be present as a result of 2-clause reduction when the residual instances include a large number ($`\mathrm{\Theta }(N)`$) of 2-clauses. As this is the case for a finite fraction of residual instances, $`G(1,1,1;T)`$ is not exponentially larger than $`L_C(T)`$. Use of the monotonicity of $`G`$ with respect to $`x_1`$ and Lemma 2 gives the announced lower bound (recognize that $`\mathrm{\Omega }(t,\alpha ,3)=t+\alpha \mathrm{log}_2\gamma (1,1;t)`$). To derive the upper bound, remark that contradictions leaves cannot be more numerous than the number of branches created through splittings; hence $`L_C(T)`$ is bounded from above by the number of splittings at smaller heights $`H`$, that is, $`{\displaystyle \underset{H<T}{}}G(0,1,1;H)`$. Once more, we use the monotonicity of $`G`$ with respect to $`x_1`$ and Lemma 2 to obtain the upper bound. The complete proof will be given in the full version.
###### Proof
(Theorem 1) By definition, a solution leaf is a node in the search tree where no clauses are left; the average number $`L_S`$ of solution leaves is thus given by $`L_S={\displaystyle \underset{H=0}{\overset{N}{}}}L(0,0,0;H)={\displaystyle \underset{H=0}{\overset{N}{}}}G(\stackrel{}{0};H)`$. A straightforward albeit useful upper bound on $`L_S`$ is obtained from $`L_SS_0(N)`$. By definition of the algorithm #DPLL, $`S_0(N)`$ is the average number of solutions of an instance with $`\alpha N`$ clauses over $`N`$ variables drawn from the random 3-SAT distribution, $`S_0(N)=2^N(7/8)^{\alpha N}`$ . This upper bound is indeed tight (to within terms that are subexponential in $`N`$), as most solution leaves have heights equal, or close to $`N`$. To show this, consider $`ϵ>0`$, and write
$$L_S\underset{H=N(1ϵ)}{\overset{N}{}}G(\stackrel{}{0};H)2^{Nϵ}\underset{H=N(1ϵ)}{\overset{N}{}}2^{NH}G(\stackrel{}{0};H)=2^{Nϵ}S_0(N)\left[1A\right]$$
with $`A=2^{Nϵ}S_0(N(1ϵ))/S_0(N)`$. From Lemma 2, $`A(\kappa +o(1))^{\alpha N}`$ with $`\kappa ={\displaystyle \frac{\gamma (0,0,1ϵ)}{7/8}}=1{\displaystyle \frac{9}{7}}ϵ^2+{\displaystyle \frac{2}{7}}ϵ^3<1`$ for small enough $`ϵ`$ (but $`\mathrm{\Theta }(1)`$ with respect to $`N`$). We conclude that $`A`$ is exponential small in $`N`$, and $`ϵ+1\alpha \mathrm{log}_2\frac{8}{7}+o(1)\frac{1}{N}\mathrm{log}_2L_S1\alpha \mathrm{log}_2\frac{8}{7}`$. Choosing arbitrarily small $`ϵ`$ allows us to establish the statement about the asymptotic behaviour of $`L_S`$ in Theorem 1.
Proposition 2, with arbitrarily small $`ϵ`$, immediately leads to Theorem 1 for $`k=3`$, for the average number of contradiction leaves, $`L_C`$, equals the sum over all heights $`T=tN`$ (with $`0t1`$) of $`L_C(N,T,\alpha )`$, and the sum is bounded from below by its largest term and, from above, by $`N`$ times this largest term. The statement on the number of leaves following Theorem 1 comes from the observation that the expected total number of leaves is $`L_S+L_C`$, and $`\omega _S(\alpha ,3)=\mathrm{\Omega }(1,\alpha ,3)\underset{t[0;1]}{\mathrm{max}}\mathrm{\Omega }(t,\alpha ,3)=\omega _C(\alpha ,3)`$. ∎
###### Proof
(Corollary 1) Let $`P_{sat}`$ be the probability that a random 3-SAT instance with $`N`$ variables and $`\alpha N`$ clauses is satisfiable. Define $`\mathrm{\#}L_{sat}`$ and $`\mathrm{\#}L_{unsat}`$ (respectively, $`L_{sat}`$ and $`L_{unsat}`$) the expected numbers of leaves in #DPLL-UC (resp. DPLL-UC) search trees for satisfiable and unsatisfiable instances respectively. All these quantities depend on $`\alpha `$ and $`N`$. As the operations of #DPLL and DPLL coincide for unsatifiable instances, we have $`\mathrm{\#}L_{unsat}=L_{unsat}`$. Conversely, $`\mathrm{\#}L_{sat}L_{sat}`$ since DPLL halts after having encountered the first solution leaf. Therefore, the difference between the average sizes #L and L of #DPLL-UC and DPLL-UC search trees satisfies $`0\mathrm{\#}LL=P_{sat}(\mathrm{\#}L_{sat}L_{sat})P_{sat}\mathrm{\#}L_{sat}`$. Hence, $`1P_{sat}\mathrm{\#}L_{sat}/\mathrm{\#}LL/\mathrm{\#}L1`$. Using $`\mathrm{\#}L_{sat}2^N`$, $`P_{sat}2^N(7/8)^{\alpha N}`$ from the first moment theorem and the asymptotic scaling for $`\mathrm{\#}L`$ given in Theorem 1, we see that the left hand side of the previous inequality tends to 1 when $`N\mathrm{}`$ and $`\alpha >\alpha _u`$. ∎
Proofs for higher values of $`k`$ are identical, and will be given in the full version.
## 4 The GUC heuristic for random SAT and COL
The above analysis of the DPLL-UC search tree can be extended to the GUC heuristic , where literals are preferentially chosen to satisfy 2-clauses (if any). The outlines of the proofs of Theorems 2 and 3 are given below; details will be found in the full version.
3-SAT. The main difference with respect to the UC case is that the two branches issued from the split are not statistically identical. In fact, the literal $`\mathrm{}`$ chosen by GUC satisfies at least one clause, while this clause is reduced to a shorter clause when $`\mathrm{}`$ is set to False. The cases $`C_21`$ and $`C_2=0`$ have also to be considered separately. With $`f_1,f_2,f_3`$ defined in the same way as in the UC case, we obtain
$`G(x_1,x_2,x_3`$ ; $`T+1)={\displaystyle \frac{1}{f_1}}G(f_1,f_2,f_3;T)+({\displaystyle \frac{1+f_1}{f_2}}{\displaystyle \frac{1}{f_1}}\left)G\right(0,f_2,f_3;T)`$ (2)
$`+`$ $`\left({\displaystyle \frac{1+f_2}{f_3}}{\displaystyle \frac{1+f_1}{f_2}}\right)G(0,0,f_3;T){\displaystyle \frac{1+f_2}{f_3}}G(0,0,0;T).`$
The asymptotic analysis of $`G`$ follows the lines of Section 3. Choosing $`f_2=f_1+f_1^2`$ i.e. $`x_1=(1+\sqrt{1+4x_2})/2+O(1/N)`$ allows us to cancel the second term on the r.h.s. of (2). Iterating relation (2), we establish the counterpart of Lemma 2 for GUC: the value of $`G`$ at height $`[tN]`$ and argument $`x_2,x_3`$ is equal to its (known) value at height 0 and argument $`y_2,y_3`$ times the product of factors $`\frac{1}{f_1}`$, up to an additive term, $`A`$, including iterates of the third and fourth terms on the right hand side of (2). $`y_2,y_3`$ are the values at ’time’ $`\tau =0`$ of the solutions of the ordinary differential equations (ODE) $`dY_2/d\tau =2m(Y_2)/(1\tau )`$, $`dY_3/d\tau =3((1+Y_2)/2Y_3)/(1\tau )`$ with ’initial’ condition $`Y_2(t)=x_2`$, $`Y_3(t)=x_3`$ (recall that function $`m`$ is defined in Theorem 2). Eliminating ’time’ between $`Y_2,Y_3`$ leads to the ODE in Theorem 2. The first term on the r.h.s. in the expression of $`\omega ^g`$ (1.2) corresponds to the logarithm of the product of factors $`\frac{1}{f_1}`$ between heights $`0`$ and $`T`$. The maximum over $`y_2`$ in expression (1.2) for $`\omega ^g`$ is equivalent to the maximum over the reduced height $`t`$ appearing in $`\omega _C`$ in Theorem 1 (see also Proposition 2). Finally, choosing $`\alpha >\alpha _u^g`$ ensures that, from the one hand, the additive term $`A`$ mentioned above is asymptotically negligible and, from the other hand, the ratio of the expected sizes of #DPLL-GUC and DPLL-GUC is asymptotically equal to unity (see proof of Corollary 1).
3-COL. The uniformity expressed by Lemma 1 holds: the subgraph resulting from the coloring of $`T`$ vertices is still Erdős-Rényi-like with edge probability $`\frac{c}{N}`$, conditioned to the numbers $`C_j`$ of vertices with $`j`$ available colors . The generating function $`G`$ of the average number of residual asignments equals $`(x_3)^N`$ at height $`T=0`$ and obeys the reccurence relation, for $`T<N`$,
$`G(x_1,x_2,x_3;T+1)`$ $`=`$ $`{\displaystyle \frac{1}{f_1}}G(f_1,f_2,f_3;T)+\left({\displaystyle \frac{2}{f_2}}{\displaystyle \frac{1}{f_1}}\right)G(0,f_2,f_3;T)`$ (3)
$`+`$ $`\left({\displaystyle \frac{3}{f_3}}{\displaystyle \frac{2}{f_2}}\right)G(0,0,f_3;T)`$
with $`f_1=(1\mu )x_1`$, $`f_2=(12\mu )x_2+2\mu x_1`$, $`f_3=(13\mu )x_3+3\mu x_2`$, and $`\mu =c/(3N)`$. Choosing $`f_1=\frac{1}{2}f_2`$ i.e. $`x_1=\frac{1}{2}x_2+O(1/N)`$ allows us to cancel the second term on the r.h.s. of (3). Iterating relation (2), we establish the counterpart of Lemma 2 for GUC: the value of $`G`$ at height $`[tN]`$ and argument $`x_2,x_3`$ is equal to its (known) value at height 0 and argument $`y_2,y_3`$ respectively, times the product of factors $`\frac{1}{f_1}`$, up to an additive term, $`A`$, including iterates of the last term in (3). An explicit calculation leads to $`G(\frac{1}{2}x_2,x_2,x_3;[tN])=e^{N\gamma ^h(x_2,x_3,t)+o(N)}+A`$ for $`x_2,x_3>0`$, where $`\gamma ^h(x_2,x_3,t)=\frac{c}{6}t^2\frac{c}{3}t+(1t)\mathrm{ln}(x_2/2)+\mathrm{ln}[3+e^{2ct/3}(2x_2/x_33)]`$. As in Proposition 2, we bound from below (respectively, above) the number of contradiction leaves in #DPLL-GUC tree by the exponential of ($`N`$ times) the value of function $`\gamma ^h`$ in $`x_2=x_3=1`$ at reduced height $`t`$ (respectively, lower than $`t`$). The maximum over $`t`$ in Theorem 3 is equivalent to the maximum over the reduced height $`t`$ appearing in $`\omega _C`$ in Theorem 1 (see also Proposition 2). Finally, we choose $`c_u^h`$ to make the additive term $`A`$ negligible. Following the notations of Corollary 1, we use $`L_{sat}3^N`$, and $`P_{sat}3^Ne^{Nc/6+o(N)}`$, the expected number of 3-colorings for random graphs from $`G(N,c/N)`$.
## 5 Conclusion and perspectives
We emphasize that the average #DPLL tree size can be calculated for even more complex heuristics e.g. making decisions based on literal degrees . This task requires, in practice, that one is able: first, to find the correct conditioning ensuring uniformity along a branch (as in the study of DPLL in the absence of backtracking); secondly, to determine the asymptotic behaviour of the associated generating function $`G`$ from the recurrence relation for $`G`$.
To some extent, the present work is an analytical implementation of an idea put forward by Knuth thirty years ago . Knuth indeed proposed to estimate the average computational effort required by a backtracking procedure through successive runs of the non–backtracking counterpart, each weighted in an appropriate way . This weight is, in the language of Section II.B, simply the probability of a branch (given the heuristic under consideration) in #DPLL search tree times $`2^S`$ where $`S`$ is the number of splits .
Since the amount of backtracking seems to have a heavy tail , the expectation is often not a good predictor in practice. Knowledge of the second moment of the search tree size would be very precious; its calculation, currently under way, requires us to treat the correlations between nodes attached to distinct branches. Calculating the second moment is a step towards the distant goal of finding the expectation of the logarithm, which probably requires a deep understanding of correlations as in the replica theory of statistical mechanics.
Last of all, #DPLL is a complete procedure for enumeration. Understanding its average-case operation will, hopefully, provide us with valuable information not only on the algorithm itself but also on random decision problems e.g. new bounds on the sat/unsat or col/uncol thresholds, or insights on the statistical properties of solutions.
Acknowledgments: The present analysis is the outcome of a work started four years ago with S. Cocco to which I am deeply indebted . I am grateful to C. Moore for numerous and illuminating discussions, as well as for a critical reading of the manuscript. I thank J. Franco for his interest and support, and the referee for pointing out Ref. , the results of which agree with the $`\alpha ^{1/(k2)}`$ asymptotic scaling of $`\omega `$ found here.
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# Carrier and Nerve Theorems in the Extension Theory
## 1. Introduction
An efficient way to investigate properties of a topological space is to divide it into pieces and examine how they are glued together. We show how to divide a general topological space and endow it with a structure that resembles a triangulation. We employ this analogy to sharpen prior results of homotopy theory, known as carrier and nerve theorems.
In order to divide a space that belongs to a class $`𝒞`$ of topological spaces we must decide what a piece is. We want it to resemble a simplex as much as possible. As a simplex is an archetype of an absolute extensor, the choice of absolute extensors for $`𝒞`$ as pieces is quite natural.
###### Definition.
A space $`Y`$ is an *absolute extensor* for a space $`X`$ if each map from a closed subset of $`X`$ into $`Y`$ extends over the entire space $`X`$. The class of absolute extensors for all spaces from a class $`𝒞`$ is denoted by $`AE(𝒞)`$. We write $`AE(X)`$ to abbreviate $`AE(\{X\})`$.
The following defines regular covers that endow a space with structures similar to triangulations. Recall that a cover is locally finite dimensional if its nerve is such.
###### Definition.
Let $`𝒞`$ be a class of topological spaces. We say that a cover is a $`𝒞`$-cover if the intersection of each non-empty collection of its elements belongs to $`𝒞`$. A locally finite $`AE(𝒞)`$-cover that is either closed and locally finite dimensional or is open is said to be *regular* for the class $`𝒞`$.
Examples of regular covers include a locally finite cover of a Euclidean space by its open balls and a cover of a finite simplicial complex by its simplices.
The main result of this paper states that the nerve of a regular cover reflects the homotopy structure of the underlying space. The result is extended to classes of spaces with bounded extension dimension, where appropriate theory of $`[L]`$-homotopy equivalences is used. Apart from generalizations, the unified approach developed in this paper delivers simplifications of proofs of earlier known theorems. In particular our carrier theorems imply nerve theorems of K. Borsuk \[3, p. 234\] and A. Weil \[15, p. 141\], while Theorem 3.4 generalizes an $`n`$-connectivity nerve theorem recently proved by A. Björner .
More specialized definitions of covers that by our definition are regular were already used in the literature. K. Kawamura gave characterizations of infinite dimensional manifolds in terms of partitions . G. de Rham used covers by convex subsets to define simple homotopy types of Riemannian manifolds . Finally nerve theorems are often used as a bridge between combinatorics and topology . Our generalization is motivated by the application of regular covers in the recent proof of characterization and rigidity theorems for Nöbeling manifolds . In A. Chigogidze conjectured that analogous characterizations hold for universal spaces for extension dimension. At the end of section 3 we give a nerve theorem for spaces of bounded extension dimension as a first step in a program to prove these conjectures.
## 2. Carrier theorems
How to extend a partial map from a subcomplex to the entire CW complex? If each map from the boundary of a Euclidean ball into the codomain extends over the ball, then the answer is easy: order cells by inclusion and construct an extension inductively. But the asphericity of the codomain (vanishing of all its homotopy groups) is a rare luxury. The same technique would work though if we were able to restrict ranges of the map on individual cells of the CW complex to aspherical subspaces of the codomain. This idea leads to the notion of a carrier and to the aspherical carrier theorem \[12, II §9\]. We generalize this notion to arbitrary spaces using a cover to replace the cell structure in the domain.
###### Definition.
A *carrier* is a function $`C:𝒢`$ from a cover $``$ of a space $`X`$ into a collection $`𝒢`$ of subsets of a topological space such that for each $`𝒜`$ if $`𝒜\mathrm{}`$, then $`_{A𝒜}C(A)\mathrm{}`$. We say that a map $`f`$ is *carried by $`C`$* if it is defined on a closed subset of $`X`$ and $`f(F)C(F)`$ for each $`F`$ (we write $`f(F)`$ for $`f(F\mathrm{dom}f)`$).
###### Carrier Theorem.
Assume that $`C:𝒢`$ is a carrier such that $``$ is a closed cover of a space $`X`$ and $`𝒢`$ is an $`AE(X)`$-cover of another space. If $``$ is locally finite and locally finitely dimensional, then each map carried by $`C`$ extends to a map of the entire space $`X`$, also carried by $`C`$.
Special cases of the Carrier Theorem follow from Michael’s Selection Theorem, as the multivalued map given by the formula $`F(x)=_{Fx}C(F)`$ is lower semi-continuous.
The proof of the Weak Carrier Theorem (later in this section) is similar to the proof of the Carrier Theorem and can be read simultaneously. The key difference is in the definition and an order of sets $`\delta _\gamma `$. The reader may compare both cases for a cover $``$ that consists of three sets with non-empty intersection.
###### Proof.
Let $`f_0`$ be a map carried by $`C`$ and let $`A=\mathrm{dom}f_0`$. Let $`\{_\gamma \}_{0<\gamma <\mathrm{\Gamma }}`$ be a transfinite sequence of all subcollections of $``$ with non-empty intersections, such that the sequence $`\{\delta _\gamma =_\gamma \}_{0<\gamma <\mathrm{\Gamma }}`$ is non-decreasing in the order by inclusion. Its existence is guaranteed by the assumption of local finite dimensionality of $``$. Let $`C_\gamma =_{E_\gamma }C(E)`$ and observe that an arbitrary map $`f`$ is carried by $`C`$ if and only if $`f(\delta _\gamma )C_\gamma `$ for each $`0<\gamma <\mathrm{\Gamma }`$. Let $`\delta _\mathrm{\Gamma }=\mathrm{}`$. We shall construct a transfinite sequence of maps $`\{f_\gamma :A_{0<\alpha \gamma }\delta _\alpha Y\}_{\gamma \mathrm{\Gamma }}`$ such that $`f_\alpha `$ extends $`f_\beta `$ for all $`0\beta \alpha \mathrm{\Gamma }`$ and $`f_\gamma (\delta _\gamma )C_\gamma `$ for each $`0<\gamma <\mathrm{\Gamma }`$. The map $`f_\mathrm{\Gamma }`$ will be an extension that we are looking for, since $`_{\gamma <\mathrm{\Gamma }}\delta _\gamma =X`$.
We proceed by transfinite induction. Fix $`\gamma \mathrm{\Gamma }`$ and assume that for each $`\alpha <\gamma `$ we already constructed $`f_\alpha `$. The map $`f_\gamma ^{}=_{\alpha <\gamma }f_\alpha `$ is well defined, continuous and its domain is closed in $`X`$ because maps $`f_\alpha `$ agree on intersections of their domains and $``$ is closed and locally finite. If $`\gamma =\mathrm{\Gamma }`$, then $`\delta _\gamma =\mathrm{}`$ and we may put $`f_\mathrm{\Gamma }=f_\mathrm{\Gamma }^{}`$. If $`\gamma <\mathrm{\Gamma }`$, then by the order of $`\delta _\gamma `$ and by inductive assumptions $`f_\gamma ^{}`$ maps $`\delta _\gamma `$ into $`C_\gamma `$. The set $`C_\gamma `$ is non-empty because $`C`$ is a carrier and is an absolute extensor for $`\delta _\gamma `$ because $`𝒢`$ is an $`AE(X)`$-cover. So $`f_\gamma ^{}`$ extends onto $`\delta _\gamma `$ to a map $`f_\gamma `$ such that $`f_\gamma (\delta _\gamma )C_\gamma `$ and our construction is finished. ∎
###### Definition.
A cover is regular for a space $`X`$ if it is regular for the class $`\{X\}`$.
###### Corollary 2.1.
If a closed cover $`𝒢`$ of a space $`Y`$ is regular for $`X\times [0,1]`$, then every two $`𝒢`$-close maps from $`X`$ into $`Y`$ are $`𝒢`$-homotopic. Moreover such homotopy exists with an additional property that if endpoints of its path lie in an element of $`𝒢`$, then the entire path lies in it.
###### Proof.
Let $`f`$ and $`g`$ denote $`𝒢`$-close maps from $`X`$ into $`Y`$. Let $`=\{F_G\}_{G𝒢}`$ be the collection of subsets of $`X\times [0,1]`$ defined by
$$F_G=(f^1(G)g^1(G))\times [0,1].$$
It is a cover of $`X\times [0,1]`$ because $`f`$ and $`g`$ are $`𝒢`$-close. Define a carrier $`C:𝒢`$ by the formula $`C(F_G)=G`$ and a map $`F:X\times \{0,1\}Y`$ by $`F(x,0)=f(x)`$ and $`F(x,1)=g(x)`$. By definition, $`F`$ is carried by $`C`$ and by the Carrier Theorem it admits an extension to the entire space $`X\times [0,1]`$, also carried by $`C`$. This extension is a $`𝒢`$-homotopy that satisfies our claim. ∎
In applications (see ), covers with slightly weaker regularity conditions are sometimes used, as in the following definition.
###### Definition.
Let $`𝒞`$ be a class of topological spaces. We say that a cover is a *weak* $`𝒞`$-cover if the union of each collection of its elements that has a non-empty intersection belongs to $`𝒞`$. A locally finite weak $`AE(𝒞)`$-cover that is either closed and locally finite dimensional or is open is said to be *weakly regular* for the class $`𝒞`$.
It follows from the Carrier Theorem that a regular cover is a weakly regular cover, but the converse is not true.
###### Definition.
Let $``$ be a cover of a space $`X`$ and let $`C:𝒢`$ be a carrier. We say that a map $`f`$ is weakly carried by $`C`$ if it is defined on a closed subset of $`X`$ and for each $`x\mathrm{dom}f`$ there exists an $`F`$ such that $`xF`$ and $`f(x)C(F)`$.
If $``$ is a cover of a space $`X`$ and $`id_{}`$ is the identity map of $``$, then a map from $`X`$ into $`X`$ is weakly carried by $`id_{}`$ if and only if it is $``$-close to the identity of $`X`$. By this observation the composition of two maps weakly carried by $`id_{}`$ does not have to be weakly carried by it.
###### Weak Carrier Theorem.
Assume that $`C:𝒢`$ is a carrier such that $``$ is an open cover of a space $`X`$ and $`𝒢`$ is a weak $`AE(X)`$-cover of another space. If $``$ is locally finite, then each map weakly carried by $`C`$ extends to a map of the entire space $`X`$, also weakly carried by $`C`$.
###### Proof.
Let $`f_0`$ be a map weakly carried by $`C`$ and let $`A=\mathrm{dom}f_0`$. Let $`\{_\gamma \}_{0<\gamma <\mathrm{\Gamma }}`$ be a transfinite sequence of all subcollections of elements of $``$ with non-empty intersections, non-decreasing in the order by inclusion. For each $`0<\gamma <\mathrm{\Gamma }`$ let $`\delta _\gamma `$ be the set of points in $`X`$ that belong exactly to all elements of $`_\gamma `$, that is, $`\delta _\gamma =_\gamma (_\gamma )`$. Let $`C_\gamma =_{E_\gamma }C(E)`$ and observe that an arbitrary map $`f`$ is weakly carried by $`C`$ if and only if $`f(\delta _\gamma )C_\gamma `$ for each $`0<\gamma <\mathrm{\Gamma }`$. Let $`\delta _\mathrm{\Gamma }=\mathrm{}`$. We shall construct a transfinite sequence of maps $`\{f_\gamma :A_{0<\alpha \gamma }\delta _\alpha Y\}_{\gamma \mathrm{\Gamma }}`$ such that $`f_\alpha `$ extends $`f_\beta `$ for all $`\beta \alpha \mathrm{\Gamma }`$ and $`f_\gamma (\delta _\gamma )C_\gamma `$ for each $`0<\gamma <\mathrm{\Gamma }`$. The map $`f_\mathrm{\Gamma }`$ will be an extension that we are looking for, since $`_{\gamma <\mathrm{\Gamma }}\delta _\gamma =X`$.
We proceed by transfinite induction. Fix $`\gamma \mathrm{\Gamma }`$ and assume that for each $`\alpha <\gamma `$ we already constructed $`f_\alpha `$. The map $`f_\gamma ^{}=_{\alpha <\gamma }f_\alpha `$ is well defined and continuous because maps $`f_\alpha `$ agree on intersections of their domains and $``$ is locally finite. By the definition of the order of sets $`\delta _\gamma `$ its domain is closed. If $`\gamma =\mathrm{\Gamma }`$, then $`\delta _\gamma =\mathrm{}`$ and we may put $`f_\mathrm{\Gamma }=f_\mathrm{\Gamma }^{}`$. If $`\gamma <\mathrm{\Gamma }`$, then by the order of $`\delta _\gamma `$ and inductive assumptions $`f_\gamma ^{}`$ maps $`\delta _\gamma `$ into $`C_\gamma `$. The set $`C_\gamma `$ is an absolute extensor for $`\delta _\gamma `$ because it is a union of elements of $`𝒢`$ that have non-empty intersection and $`𝒢`$ is a weak $`AE(X)`$-cover. So $`f_\gamma ^{}`$ extends onto $`\delta _\gamma `$ to a map $`f_\gamma `$ such that $`f_\gamma (\delta _\gamma )C_\gamma `$ and our construction is finished. ∎
###### Corollary 2.2.
If an open cover $`𝒢`$ of a space $`Y`$ is weakly regular for $`X\times [0,1]`$, then every two $`𝒢`$-close maps from $`X`$ into $`Y`$ are $`\mathrm{st}𝒢`$-homotopic, where $`\mathrm{st}𝒢`$ denotes the star of $`𝒢`$. Moreover such homotopy exists with an additional property that if endpoints of its path lie in an element of $`𝒢`$, then the entire path lies in its star.
We omit the proof as it is similar to the proof of Corollary 2.1.
###### Remark.
In the Weak Carrier Theorem the assumption that $``$ is locally finite may be omitted if $`𝒢`$ is open and $`X`$ is paracompact. To prove this, it suffices to find a locally finite open refinement $``$ of $``$ and a carrier $`D:`$ such that $`f`$ is weakly carried by $`CD`$ and $`HD(H)`$ for each $`H`$. Then, by the Weak Carrier Theorem, $`f`$ extends over $`X`$ to a map weakly carried by $`CD`$, which is obviously weakly carried by $`C`$. Let $``$ be any locally finite refinement of an open cover $`\{Ff^1(C(F)):F\}`$ and let $`D:`$ be any map such that $`HD(H)f^1(C(D(H)))`$ for each $`H`$. Then $`f`$ is carried by $`CD`$ because for each $`xX`$ there exists $`H`$ such that $`xHD(H)f^1(C(D(H)))`$ so $`f(x)(CD)(H)`$. By the definition $`f`$ is weakly carried by $`CD`$ and the proof is finished.
Similarly in Corollary 2.2 the assumption of local finiteness of $`𝒢`$ may be replaced by paracompactness of $`X`$.
## 3. Nerve theorems
Nerve theorems give conditions under which the nerve of a cover is equivalent to the underlying space. First examples of such theorems are attributed to K. Borsuk \[3, p. 234\] (for closed covers) and A. Weil \[15, p. 141\] (for open covers), both for homotopy equivalences. Since then, several generalizations were made. First generalizations by W. Holsztyński and J. N. Haimov relaxed conditions on the cover. Next weak homotopy equivalences were studied in this context by M. McCord and weak $`n`$-homotopy equivalences by A. Björner .
Let $`K`$ denote an arbitrary simplicial complex. An open star $`\mathrm{st}v`$ of a vertex $`vK`$ is the complement of the union of all simplices of $`K`$ that do not contain $`v`$. A barycentric star $`\mathrm{bst}v`$ of a vertex $`vK`$ is the union of all simplices of the first barycentric subdivision of $`K`$ that contain $`v`$.
###### Lemma 3.1.
If the cover of a simplicial complex (with the metric topology) by the collection of open (or barycentric) stars of its vertices is locally finite (and locally finite dimensional), then it is regular for the class of metric spaces. If additionally the complex is locally countable, then the cover is regular for the class of normal spaces.
By \[10, Theorem 11.7, p. 109\] and by \[10, Theorem 7.1, p. 43\], it suffices to show that each non-empty intersection of a collection of elements of the cover is contractible. We omit details of the proof, which is based on the observation that every point of the intersection is connected by a line with the barycenter of centers of stars that are intersected.
We shall use the following notation.
###### Definition.
Let $`𝒮_{}`$ denote the collection of open stars of vertices of the nerve $`N()`$ of a point-finite cover $``$. Let $`_{}`$ denote the collection of its barycentric stars.
By Lemma 3.1, covers $`𝒮_{}`$ and $`_{}`$ are regular for the class of normal spaces. Their structure mirrors the structure of $``$, in the sense of the following definition.
###### Definition.
We say that covers are *isomorphic* if there exists a carrier from one of them onto the other, which is *invertible* (i.e., is bijective and its inverse is a carrier).
Observe that if the codomain of a carrier $`C`$ is equal to the domain of a carrier $`D`$, then a composition of two maps carried respectively by $`C`$ and $`D`$ is carried by $`DC`$.
###### Lemma 3.2.
Let $``$ be a point-finite cover and let $`v(F)`$ denote the vertex of $`N()`$ corresponding to a set $`F`$. The functions $`S:𝒮_{}`$ and $`B:_{}`$ defined by $`S(F)=\mathrm{st}v(F)`$ and $`B(F)=\mathrm{bst}v(F)`$ are invertible carriers.
###### Proof.
Obviously $`S`$ and $`B`$ are bijections. We shall prove that $`B`$ and $`S^1`$ are carriers. Then, since a function $`I:_{}𝒮_{}`$ defined by the formula $`I(\mathrm{bst}v)=\mathrm{st}v`$ is a carrier and $`S^1IB=id_{}`$, functions $`B^1=S^1I`$ and $`S=IB`$ are carriers too.
To prove that $`B`$ is a carrier, assume that $`\{F_i\}`$ is a collection of elements of $``$ with non-empty intersection. Then the nerve $`N()`$ contains a simplex $`\sigma `$ spanned by vertices $`\{v(F_i)\}`$ and each barycentric star $`\mathrm{bst}v(F_i)`$ contains the barycenter of $`\sigma `$. Therefore the intersection of $`\{B(F_i)\}`$ is non-empty.
To prove that $`S^1`$ is a carrier, assume that $`\{\mathrm{st}v(F_i)\}`$ is a collection of elements of $`𝒮_{}`$ with non-empty intersection. As an open star of a vertex $`v(F_i)`$ contains only points of simplices that contain $`v(F_i)`$, then the intersection of the set $`\{\mathrm{st}v(F_i)\}`$ contains only points of simplices that contain all vertices $`v(F_i)`$. So the nerve $`N()`$ contains a simplex spanned by $`\{v(F_i)\}`$ and the intersection of $`\{F_i\}`$ is non-empty. ∎
Putting everything together we obtain nerve theorems for homotopy equivalences. We state the theorem for closed covers, which generalizes a nerve theorem by J. N. Haimov . An analogous theorem for open covers may also be proved. We do not state it here, as it turns out to be equivalent to the nerve theorem by A. Weil .
###### Theorem 3.3.
Assume that a closed cover $``$ of a normal space $`X`$ is regular for the class of metric spaces. If $``$ is star-countable, then $`X`$ and the nerve of $``$ are homotopy equivalent.
The main theorem of states the same conclusion under the additional assumption that $``$ is star-finite and $`X`$ is paracompact.
###### Proof.
Let $`B:_{}`$ be an invertible carrier as defined in Lemma 3.2. By the Carrier Theorem and Lemma 3.1 there exist $`f:XN()`$ and $`g:N()X`$ carried by $`B`$ and $`B^1`$ respectively. Then $`gf`$ is carried by $`B^1B`$ so it is $``$-close to $`id_X`$ and by Corollary 2.1 $`g`$ is a homotopy inverse of $`f`$. Analogously $`f`$ is a homotopy inverse of $`g`$ so $`X`$ and $`N()`$ are homotopy equivalent. ∎
We turn our attention to spaces with bounded dimension. First we prove a nerve theorem for the class of at most $`n`$-dimensional spaces. Next the case of extension dimension is studied. To avoid anomalies of the dimension of general topological spaces from now on all spaces are assumed to be separable metric.
###### Theorem 3.4.
If $``$ is an open cover of a space $`X`$, weakly regular for the class of at most $`n`$-dimensional spaces, then each canonical map $`\varkappa :XN()`$ induces isomorphisms on homotopy groups of dimensions less than $`n`$.
A canonical map into the nerve of a cover is a map induced by a partition of unity subordinated to this cover, by interpreting its values as barycentric coordinates of points in the nerve. Kuratowski’s $`\varkappa `$-map \[7, p. 321\] is an example of such a map.
It follows from the Excision Theorem that an open cover $``$ of a locally $`(n1)`$-connected space is weakly regular for the class of at most $`n`$-dimensional spaces if and only if an intersection of each collection $`𝒜`$ is $`(n\left|𝒜\right|)`$-connected. This relates our theorem to Björner’s nerve theorem .
###### Proof.
For each subcollection $`𝒜`$ of $``$ with non-empty intersection let $`v(𝒜)`$ denote the simplex of the nerve $`N()`$ spanned by vertices $`\{v(A)\}_{A𝒜}`$.
Let $`\lambda :N()^{(n)}X`$ denote a map from the $`n`$-dimensional skeleton of $`N()`$ into $`X`$, weakly carried by the carrier $`S^1`$ defined in Lemma 3.2. Fix $`k<n`$. For each map $`\phi :S^kN()`$ pick a map $`\phi ^{}:S^kN()^{(n)}`$, homotopic to $`\phi `$, such that if $`\phi (x)\sigma `$, then $`\phi ^{}(x)\sigma ^{(n)}`$ for each simplex $`\sigma `$ of $`N()`$. It exists by the Cellular Approximation Theorem.
To prove that $`\varkappa `$ induces an epimorphism on $`k`$th homotopy groups we shall observe that for each $`\phi :S^kN()`$ the map $`\varkappa (\lambda \phi ^{})`$ is $`𝒮_{}`$-close to $`\phi `$, so by Corollary 2.2 they are homotopic. Fix $`xS^k`$. From the definition of a weakly carried map there exists $`F`$ such that $`\phi ^{}(x)\mathrm{st}v(F)`$ and $`\lambda (\phi ^{}(x))F`$. From the definition of $`\phi ^{}`$ we have $`\phi (x)\mathrm{st}v(F)`$. From the definition of $`\varkappa `$ we have $`\varkappa (F)\mathrm{st}v(F)`$. Then $`\varkappa (\lambda (\phi ^{}(x)))\mathrm{st}v(F)`$ and we are done.
To prove that $`\varkappa `$ induces a monomorphism on $`k`$th homotopy groups we shall observe that for each $`\psi :S^kX`$ the map $`\lambda ((\varkappa \psi )^{})`$ is $``$-close to $`\psi `$, so by Corollary 2.2 they are homotopic, as dimension of $`S^k\times [0,1]`$ is at most $`n`$. Fix $`xS^k`$ and let $`𝒜=\{F:\psi (x)F\}`$. Then $`(\varkappa \psi )(x)v(𝒜)`$ so $`(\varkappa \psi )^{}(x)v(𝒜)^{(n)}`$. By the definition of $`𝒮_{}`$, $`v(𝒜)\mathrm{st}v(F)\mathrm{}`$ if and only if $`F𝒜`$, so $`\lambda (\varkappa \psi )^{}(x)F`$ for some $`F𝒜`$. Therefore if $`\varkappa \psi `$ is null-homotopic, then $`\lambda ((\varkappa \psi )^{})`$ and so is $`\psi `$. ∎
###### Remark.
Let $`\sigma (𝒜)=\mathrm{Cl}(𝒜)(𝒜)`$ be the closure of the set of points in $`X`$ that belong exactly to elements of $`𝒜`$. Then the map $`K(\sigma (𝒜))=v(𝒜)`$, where $`𝒜`$ runs over all subcollections of $``$ that have non-empty intersections, is a carrier. To prove that, observe that if $`x_i\sigma (𝒜_i)`$ then $`𝒜=\{F:xF\}𝒜_i`$ for each $`i`$, so $`_iv(𝒜_i)v(𝒜)\mathrm{}`$. Every canonical map into the nerve of a cover is carried by the carrier $`K`$.
We finish with a generalization of Theorem 3.4 to a class of spaces with uniformly bounded extension dimension. The survey is a good source of information about notions of the extension dimension and the $`[L]`$-homotopy; here we will only recall basic definitions.
###### Definition.
Let $`L`$ be an arbitrary CW complex. A space $`X`$ has extension dimension less than or equal to $`[L]`$ if $`L`$ is an absolute extensor for $`X`$. The class of absolute extensors for at most $`[L]`$-dimensional spaces is denoted by $`AE[L]`$.
We say that maps $`f,g:XY`$ are $`[L]`$-homotopic if for each at most $`[L]`$-dimensional space $`Z`$, each pair $`A`$, $`B`$ of disjoint closed subsets of $`Z`$ and each map $`h:ZX`$ there exists an extension of a map $`fh_{|A}gh_{|B}`$ to the entire space $`Z`$.
We shall need the following analogue of Corollary 2.2.
###### Lemma 3.5.
If $`𝒢`$ is a locally finite closed weak $`AE[L]`$-cover of a separable metric space, then every two $`𝒢`$-close maps are $`[L]`$-homotopic.
###### Proof.
Name the maps considered by $`f,g:XY`$, where $`Y`$ denotes the space covered by $`𝒢`$. We are going to prove that $`f`$ and $`g`$ are $`[L]`$-homotopic directly from the definition. Let $`Z`$ and $`h`$ be as in the definition of the $`[L]`$-homotopy.
Let $`=\{E_G\}_{G𝒢}`$ be the collection of subsets of $`Z`$ defined by
$$E_G=h^1(f^1(G)g^1(G)).$$
It is a cover of $`Z`$ because $`f`$ and $`g`$ are $`𝒢`$-close. The map $`C:𝒢`$ given by the formula $`C(F_G)=G`$ carries $``$ into $`𝒢`$. By the definition $`H=fh_{|A}gh_{|B}`$ is carried by $`C`$.
Since $`𝒢`$ is a locally finite, closed cover, so is $``$. Then by \[6, p. 393\], since $`Z`$ is separable metric, there is a locally finite open cover $`=\{F_E\}_E`$ of $`Z`$ such that $`EF_E`$ for each $`E`$ and the function $`F_EE`$ is an invertible carrier of $``$ onto $``$. Let $`D:`$ be a function such that $`FD(F)`$ for each $`F`$. Then the identity map $`id_Z`$ is weakly carried by $`D`$, and $`Hid_Z`$ is weakly carried by $`CD`$. By the Weak Carrier Theorem, the map $`H`$ extends to the entire space $`Z`$ and the proof is finished. ∎
The following is a nerve theorem for $`[L]`$-homotopy.
###### Definition.
We say that topological spaces $`X`$ and $`Y`$ are $`[L]`$-homotopy equivalent if there exist mappings $`f:XY`$ and $`g:YX`$ such that $`gf`$ and $`fg`$ are $`[L]`$-homotopic to $`id_X`$ and $`id_Y`$ respectively.
###### Theorem 3.6.
If two at most $`[L]`$-dimensional spaces have isomorphic closed covers, regular for the class of at most $`[L]`$-dimensional spaces, then they are $`[L]`$-homotopy equivalent.
We omit the proof, which is a straightforward application of the Carrier Theorem and Lemma 3.5, similar to the proof of Theorem 3.3.
###### Remark.
Observe that an $`[S^{n1}]`$-homotopy is J. H. C. Whitehead’s $`n`$-homotopy, and by Whitehead’s characterization a map between locally $`(n1)`$-connected $`n`$-dimensional spaces is an $`[S^{n1}]`$-homotopy equivalence if and only if it induces isomorphisms on homotopy groups of dimensions less than $`n`$. This connects theorem 3.4 with Theorem 3.6.
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# The Color–Flavor Transformation of induced QCD
## 1 Introduction
The complexity of quantum chromodynamics (QCD) originates from the random character of the gauge field in the low-energy regime, while at high energy (small scales) the theory is asymptotically free. One of the most successful approaches to analyse QCD in the non-perturbative domain is the lattice formulation due to Wilson , where the strong coupling regime becomes natural. On the other hand, the continuum limit of lattice QCD corresponds to the weakly-coupled regime. It was shown by Gross and Witten that two-dimensional lattice QCD can be solved exactly in the large-$`N_c`$ limit and there is a third-order weak- to strong-coupling phase transition. The problem, which persists over the years is, if such a phase transition occurs in the realistic $`SU(3)`$ four-dimensional gauge theory.
One of the ways to attack this problem is to use some form of *duality* which appears in the lattice theory. The general idea of duality is that a given theory can have two, or more equivalent formulations, with different sets of fundamental variables. Usually these dual formulations are related by interchanging a parameter, e.g. the electromagnetic coupling constant $`e^2`$, with its inverse $`1/e^2`$. For example, in the weak-coupling limit, the action of the Abelian lattice gauge model can be approximated by the Villain form (see e.g. ) which allows to define the dual variables. Similarly, it was shown that there is a duality transformation from the compact $`U(1)`$ gauge theory into a non-compact Abelian Higgs model .
The dual variables are in general not defined on the original lattice but on a *dual lattice*. For a hypercubic lattice, it is obtained by shifting the original lattice by half of the lattice spacing in all dimensions. Thus, the lattice duality not only transforms the variables of the functional integral, but also incorporates a transfer to the dual lattice.
In the case of a nonabelian lattice gauge theory, one can reformulate the model in terms of plaquette variables . Recently, following the idea of , Diakonov and Petrov applied a Fourier transformation to write down the Jacobian from the link variables to the plaquette variables in $`d=3`$ gluodynamics for the gauge group SU$`(2)`$; *in fine* the dual lattice is composed of tetrahedra representing $`6j`$-symbols, with links of arbitrary lengths. In the continuum limit this effective theory is equivalent with quantum gravity with the Einstein-Hilbert action. Unfortunately, this scheme seems difficult to apply to higher-rank gauge groups and in $`d>3`$ dimensions.
Another approach to analyze non-perturbative QCD, the so-called “Induced QCD”, was initiated in the 90s . This direction originated from A.D. Sakharov’s idea to treat the gauge theory as induced quantum field theory. In the Kazakov-Migdal model the Wilson lattice action is induced by the auxiliary heavy scalar matrix fields which are taken in an adjoint representation of the gauge group $`SU(N_c)`$. This field can be diagonalized by a gauge transformation so that in the large $`N_c`$ limit the functional integral over eigenvalues of the matrix field, which serves as a master field, is saturated by a saddle point of the effective action. The induced action obtained via integration over the auxiliary fields contains the traces of products of the link variables along all possible contours. However in the large mass limit, only one-plaquette loops survive and the Wilson action is recovered.
The interest in the Kazakov-Migdal model was mainly inspired by its exact solvability in the large $`N_c`$ limit but its continuum limit is questionable. Moreover, for inducing fields in adjoint representation there is an extra local $`Z_N`$ symmetry which leads to infinite string tension (the so-called local confinement) rather than the conventional area law for the Wilson loop . That is also the case for the adjoint fermion model of the induced QCD .
Another example of the QCD inducing model was proposed by Bander and Hamber . In this approach the Wilson action is recovered if the number of flavors of the auxiliary fields goes to infinity simultaneously with the mass. This model does not suffer from the extra gauge symmetry but is not solvable even in $`d=2`$.
In the present note, based on a work in collaboration with S. Nonnenmacher , I discuss a different approach to treat a similar type of inducing model. Our construction starts from an inducing theory similar with , already introduced in , and applies a certain duality transformation, namely the “color–flavor transformation” . We have recently applied this transformation to the lattice SU$`(N_c)`$ model in the strong-coupling limit, which describes quarks coupled with a background gauge field , see also paper . Schlittgen and Wettig also independently applied the SU$`(N)`$ color-flavor transformation to a similar, yet different QCD-inducing model . A very interesting formalism to induce lattice gauge model is suggested by Budczies and Zirnbauer in paper , where instead of heavy inducing fields a finite number $`N_b`$ of auxiliary boson flavors was coupled to the gauge field. In this framework $`U(N_c)`$ gauge theory is induced when $`N_b`$ exceeds $`N_c`$ and the boson mass is lowered to a critical point. In the present note I would like to sum up our investigation of a possibility to apply another “color–flavor motivated” approach, which is related with some modification of the model by Bander and Hamber.
## 2 A model of induced lattice gauge theory
### 2.1 Wilson’s lattice action
We consider a Euclidean U$`(N_c)`$ pure gauge action (no quarks) in $`d`$ dimensions, placed on a hypercubic lattice with lattice constant $`a`$. The choice of U$`(N_c)`$ instead of the realistic SU$`(N_c)`$ highly simplifies the subsequent color-flavor transformation.
The lattice sites are labeled by integer vectors $`n=(1,\mathrm{},n_d)`$, the gauge matrix variables
$$U_\mu (n)U_{n,n+\mu }U(n+{\scriptscriptstyle \frac{\widehat{\mu }}{2}})=\mathrm{exp}\left(iagA_\mu (na+\frac{a\widehat{\mu }}{2})\right)\mathrm{U}(N_c)$$
(1)
are placed on the lattice links $`n+\widehat{\mu }/2`$ (we label links by their *middle points*), leaving the site $`n`$ in any of the “positive” directions $`\mu =1,\mathrm{},d`$. The plaquettes are either labeled by an independent index $`p`$, or by triplets of the form $`(n,\pm \mu ,\pm \nu )`$. For instance, the plaquette $`(n,\mu ,\nu )`$ contains the links $`n+\widehat{\mu }/2`$ and $`n+\widehat{\nu }/2`$. To fix an ordering between the directions, we will in general assume that $`1\mu <\nu d`$. Notice that the same plaquette corresponds to the triplets $`(n,\mu ,\nu )`$ and $`(n+\widehat{\mu },\mu ,\nu )`$ (as well as two other triplets).
The Wilson pure gauge action is given by a sum over all elementary plaquettes:
$$S_{\mathrm{gluons}}=\beta _W\underset{p}{}\mathrm{Tr}\left(U_P(p)+U_P^{}(p)\right).$$
(2)
Here $`\beta _W`$ is the lattice coupling, which is related to the bare continuum coupling constant $`g`$ through
$$\beta _W=\frac{a^{d4}}{2g^2}.$$
(3)
The plaquette field $`U_P`$ is defined as an ordered product of the link variables along the boundary of the given plaquette:
$$U_P(n,\mu ,\nu )=U(n+{\scriptscriptstyle \frac{\widehat{\mu }}{2}})U(n+\widehat{\mu }+{\scriptscriptstyle \frac{\widehat{\nu }}{2}})U^1(n+\widehat{\nu }+{\scriptscriptstyle \frac{\widehat{\mu }}{2}})U^1(n+{\scriptscriptstyle \frac{\widehat{\nu }}{2}}).$$
(4)
The partition function is defined as
$$Z=𝒟Ue^{S_{\mathrm{gluons}}},$$
(5)
the invariant measure of integration is defined as a product over all links $`𝒟U=\underset{n,\mu }{}dU(n+{\scriptscriptstyle \frac{\widehat{\mu }}{2}})`$ and $`dU(n+{\scriptscriptstyle \frac{\widehat{\mu }}{2}})`$ is the Haar measure on the group U$`(N_c)`$.
Using a generalized Baker-Campbell-Hausdorff formula, one can relate the continuum field strength tensor
$$F_{\mu \nu }=_\mu A_\nu _\nu A_\mu i[A_\mu ,A_\nu ]$$
with the plaquette matrices as follows:
$$U_P(n,\mu ,\nu )=e^{iga^2F_{\mu \nu }(n)+O(a^3)}.$$
(6)
Expansion in the lattice spacing up to the second order then yields
$$\mathrm{Tr}U_P(n,\mu ,\nu )N_c+iga^2\mathrm{Tr}F_{\mu \nu }\frac{g^2a^4}{2}\mathrm{Tr}F_{\mu \nu }^2,$$
so that the standard Yang-Mills action is recovered in the continuum limit:
$$S_{cont}=\frac{1}{2}d^dx\mathrm{Tr}F_{\mu \nu }^2.$$
(7)
Wilson’s lattice gauge action (2) is written via the plaquette matrices $`U_P(p)`$ while the variables of the integration measure are the link matrices $`U(n+{\scriptscriptstyle \frac{\widehat{\mu }}{2}})`$. It is possible to explicitly transfer the integration from link to plaquette matrices , but this procedure is technically involved and can be applied only in a few simple cases. We would like to investigate a possibility to apply another approach, which starts from a modification of the QCD inducing model of Bander and Hamber .
### 2.2 Description of the model and its induced gauge action
Let us consider a massive complex bosonic field $`\varphi (n)`$ placed on the lattice sites $`n`$. This field has “flavor”components $`\varphi ^{(\pm \mu ,\pm \nu )}(n)`$ associated to each of the $`2d(d1)`$ plaquette $`\{(n,\pm \mu ,\pm \nu );1\mu <\nu d\}`$ adjacent to the site $`n`$ (see Figures 1,3). The field furthermore decomposes into two “chiral components” $`\varphi _R^{(\mu ,\nu )}(n)`$ and $`\varphi _L^{(\mu ,\nu )}(n)`$, which are hopping in opposite directions. These fields will be referred to as the “left” component and the “right” component respectively.
The chiral bosonic field can be thought of as an $`N_b`$-component vector in an auxiliary “flavor” space (here the index $`b`$ stands for “bosonic”). The number of “flavors” has to be a multiple of $`2d(d1)`$, that is the dimension of the “flavor” space is $`N_b=n_b\times 2d(d1)`$, where $`n_b^{}`$ is the number of “generations” of the bosonic field. The bosonic field $`\varphi `$ also transforms as a vector through the gauge group $`U(N_c)`$, so it contains “color indices” $`i=1,\mathrm{},N_c`$ besides the “flavor” indices $`a=1,\mathrm{},N_b`$. All flavor components have the same mass $`m_b`$.
This model (already considered in \[25, Chap. 5\] and ) is different from the model usually considered in induced QCD , where the flavor degrees of freedom of the auxiliary fields are not associated with plaquettes. As we will show below, this structure will induce a Wilson-type action in a cleaner way than in the previous models.
To complete our notations, we shall fix the orientation of the plaquettes, in order to define the hopping and the “left” and the “right” components in all $`2`$-dimensional planes on the lattice. On any plane $`(\mu ,\nu )`$ with $`1\mu <\nu d`$ the ‘left’ chiral component hops on the plaquette $`(n,\mu ,\nu )`$ as
$$\varphi _L(n)\varphi _L(n+\widehat{\mu })\varphi _L(n+\widehat{\mu }+\widehat{\nu })\varphi _L(n+\widehat{\nu })\varphi _L(n),$$
(8)
and the ‘right’ chiral component hops in the opposite direction (see Fig. 2). This practically means that the “kinetic” part of the action contains the term (cf. Eq. (10)):
$$\varphi _L^{(\mu ,\nu )}(n+\widehat{\mu })U(n+{\scriptscriptstyle \frac{\widehat{\mu }}{2}})\varphi _L^{(\mu ,\nu )}(n).$$
Let us now group the fields surrounding a given plaquette $`p=(n,\mu ,\nu )`$ into the following *plaquette quadruplets*:
$$\varphi _L(p)\stackrel{\mathrm{def}}{=}\left(\begin{array}{c}\varphi _L^{(\mu ,\nu )}(n)\\ \varphi _L^{(\mu ,\nu )}(n+\widehat{\mu })\\ \varphi _L^{(\mu ,\nu )}(n+\widehat{\mu }+\widehat{\nu })\\ \varphi _L^{(\mu ,\nu )}(n+\widehat{\nu })\end{array}\right);\varphi _R(p)\stackrel{\mathrm{def}}{=}\left(\begin{array}{c}\varphi _R^{(\mu ,\nu )}(n)\\ \varphi _R^{(\mu ,\nu )}(n+\widehat{\mu })\\ \varphi _R^{(\mu ,\nu )}(n+\widehat{\mu }+\widehat{\nu })\\ \varphi _R^{(\mu ,\nu )}(n+\widehat{\nu })\end{array}\right)$$
(9)
Then the plaquette action of the “left” bosonic massive field may be written in the concise form
$$S_L(p)=\varphi _L^{}(p)M_L(p)\varphi _L(p)$$
(10)
with the $`4N_c\times 4N_c`$ matrix
$$M_L(n,\mu ,\nu )\stackrel{\mathrm{def}}{=}\left(\begin{array}{cccc}m_b& 0& 0& U^1(n+{\scriptscriptstyle \frac{\widehat{\nu }}{2}})\\ U(n+{\scriptscriptstyle \frac{\widehat{\mu }}{2}})& m_b& 0& 0\\ 0& U(n+\widehat{\mu }+{\scriptscriptstyle \frac{\widehat{\nu }}{2}})& m_b& 0\\ 0& 0& U^1(n+\widehat{\nu }+{\scriptscriptstyle \frac{\widehat{\mu }}{2}})& m_b\end{array}\right)$$
(11)
Similarly, “right” bosonic action associated with the plaquette $`p`$ reads:
$$S_R(p)=\varphi _R^{}(p)M_L^{}(p)\varphi _R(p).$$
(12)
We define the full partition function as:
$$𝒵=𝒟\overline{\varphi }_{L,R}𝒟\varphi _{L,R}𝒟U\mathrm{exp}\left\{\underset{p}{}S_L(p)+S_R(p)\right\}.$$
(13)
### 2.3 Integration over auxiliary fields
We now show that the model (13) induces a lattice gluodynamics which reduces to Wilson’s pure gauge action for suitably chosen parameters. We can treat plaquettes separately, since each field component is associated with one and only plaquette. The integration over the ‘left’ auxiliary fields in (10) yields
$$d\overline{\varphi }_L(p)d\varphi _L(p)\mathrm{exp}\{S_L(p)\}=\mathrm{Det}\left(M_L(p)\right)^{n_b}=\mathrm{Det}(m_b^4U_P(p))^{n_b},$$
(14)
where $`n_b`$ is the number of “generations” of the auxiliary fields. The integration over the “right” components is similar, with $`U_P`$ replaced by $`U_P^{}`$. Thus, the integration over the auxiliary bosonic fields exactly yields the pure gauge effective action
$$\begin{array}{cc}\hfill S_{\mathrm{plaq}}& =n_b\underset{p}{}[\mathrm{ln}\mathrm{Det}(1\beta _bU_P(p))+\mathrm{ln}\mathrm{Det}(1\beta _bU_P^{}(p))]\hfill \\ & =n_b\underset{p}{}\mathrm{Tr}[\mathrm{ln}(1\beta _bU_P(p))+\mathrm{ln}(1\beta _bU_P^{}(p))],\hfill \end{array}$$
where we skipped a mass-dependent prefactor, and set $`\beta _b=m_b^4`$. As opposed to the former inducing models , this action does not contain any term related with larger loops. Expanding this action for small parameter $`\beta _b`$ (that is, large $`m_b`$), we get
$$S_{\mathrm{plaq}}=n_b\beta _b\mathrm{Tr}\underset{p}{}\left(U_P(p)+U_P^{}(p)\right)+𝒪(n_b\beta _b^2).$$
(15)
This coincides with Wilson’s action (2) if we identify
$$\beta _W=n_b\beta _bg^2=\frac{a^{d4}}{2n_b\beta _b}=\frac{a^{d4}m_b^4}{2n_b}.$$
(16)
#### Remarks on the continuum limit
Let us say a few words about the continuum limit of the model. In $`d=2`$ and $`d=3`$, the physical coupling constant $`g`$ can remain fixed as one lets the lattice spacing $`a`$ go to $`0`$ and simultaneously $`\beta _W=n_b\beta _b\mathrm{}`$. In $`d=4`$, this limit corresponds to the asymptotically free continuum theory. Recall that the mass of the auxiliary fields is measured in units of the lattice spacing: $`m_b=m\times a`$. Therefore the auxiliary field becomes non-observable in the continuum limit if the corresponding correlation length $`\xi =(am)^1=\beta _b^{1/4}`$ stays finite. We assumes it remains small enough to justify the expansion (15).
We end this section with one more comment. It is possible to consider a fermionic counterpart of the bosonic action (13), which induces a similar pure gauge effective action. One has to replace the bosonic auxiliary fields by fermionic (anticommuting) fields $`\psi _{L,R}`$, $`\overline{\psi }_{L,R}`$ which carry color indices and flavor indices related to plaquettes, exactly as for the multiplets (9); one can consider $`n_f`$ “generations” of these fermions. After integrating over them, one obtains the effective pure gauge action
$$S=n_f\underset{p}{}\mathrm{Tr}[\mathrm{ln}(1+\beta _fU_P(p))+\mathrm{ln}(1+\beta _fU_P^{}(p))].$$
(17)
Here $`\beta _f=m_f^4`$, where $`m_f`$ is the fermion mass. Clearly, we once more recover the conventional Wilson action (2) for small values of $`\beta _f`$.
Having considered both bosonic and fermionic induced lattice gluodynamics, we can also represent the effective action as the ratio of “fermionic” and “bosonic” determinants:
$$\mathrm{exp}\left[\frac{a^{d4}}{2g^2}\mathrm{Tr}(U_P+U_P^{})\right]\frac{\left[\mathrm{Det}(1+\beta _fU_P)\mathrm{Det}(1+\beta _fU_P^{})\right]^{n_f}}{\left[\mathrm{Det}(1\beta _bU_P)\mathrm{Det}(1\beta _bU_P^{})\right]^{n_b}}.$$
(18)
Thus, for small couplings $`\beta _b`$, $`\beta _f`$, the partition function can be represented by the following superintegral
$$𝒵=𝒟U𝒟\psi 𝒟\overline{\psi }\mathrm{exp}\left[\overline{\psi }_{L,a}^i(\delta ^{ij}+\beta _L(U_P)^{ij})\psi _{L,a}^j\overline{\psi }_{R,a}^i(\delta ^{ij}+\beta _R(U_P^{})^{ij})\psi _{R,a}^j\right].$$
(19)
The composite field $`\psi ,\overline{\psi }`$ includes both bosonic and fermionic variables, which are distinguished by the “flavor” index $`a`$.
There are therefore several ways to induce Wilson’s lattice gauge action. In all cases, the action (2) with fixed $`\beta _W`$ can be recovered in the limit of large mass and large number of generations of the auxiliary fields (cf. Eq. (16)). Below we will restrict our considerations to the bosonic model (13).
## 3 Color-flavor transformation of the inducing theory
Though the equivalence between the model (13) and Wilson’s gluodynamics can be established only in the limit of large number of generations of the auxiliary field, we would like to study some properties of the underlying theory with a single bosonic “generation” $`(n_b=1)`$, that is with a flavor space of dimension $`N_b=2d(d1)`$. This assumption simplifies the application of the color-flavor transformation .
Let us consider the interaction term of the bosonic action (13) on a given link $`(n+\widehat{\mu }/2)`$ of the $`d`$-dimensional hypercubic lattice. The path ordered product of the link matrices defined by the “left” action (10) is depicted in Fig. 3. There are $`2d2`$ plaquettes which share this common link. Above we have used the plaquette quadruplets (9) in order to write the plaquette action concisely. Now we will rather decompose the full action into a sum over *links*, which forces us to gather the auxiliary bosonic fields into two series of *site-link multiplets*, in order to include all fields coupled by $`U(n+{\scriptscriptstyle \frac{\widehat{\mu }}{2}})`$ or $`U^{}(n+{\scriptscriptstyle \frac{\widehat{\mu }}{2}})`$. We thus define two chirally-conjugated multiplets associated with the site $`n`$ and the link $`(n+\widehat{\mu }/2)`$:
$$\mathrm{\Psi }(n;\mu )\stackrel{\mathrm{def}}{=}\left(\begin{array}{c}\varphi _R^{(\mu ,\mu +1)}(n)\\ \varphi _L^{(\mu ,\mu 1)}(n)\\ \mathrm{}\\ \varphi _R^{(\mu ,d)}(n)\\ \varphi _L^{(\mu ,d)}(n)\\ \varphi _R^{(\mu ,1)}(n)\\ \varphi _L^{(\mu ,1)}(n)\\ \mathrm{}\\ \varphi _R^{(\mu ,\mu +1)}(n)\\ \varphi _L^{(\mu ,\mu 1)}(n)\end{array}\right);\mathrm{\Phi }(n;\mu )\stackrel{\mathrm{def}}{=}\left(\begin{array}{c}\varphi _L^{(\mu ,\mu +1)}(n)\\ \varphi _R^{(\mu ,\mu 1)}(n)\\ \mathrm{}\\ \varphi _L^{(\mu ,d)}(n)\\ \varphi _R^{(\mu ,d)}(n)\\ \varphi _L^{(\mu ,1)}(n)\\ \varphi _R^{(\mu ,1)}(n)\\ \mathrm{}\\ \varphi _L^{(\mu ,\mu +1)}(n)\\ \varphi _R^{(\mu ,\mu 1)}(n)\end{array}\right)$$
(20)
The multiplets $`\mathrm{\Psi }(n;\mu )`$ and $`\mathrm{\Phi }(n;\mu )`$ associated with the site $`n`$ and link $`n\widehat{\mu }/2`$ are obtained from the ones above by flipping the *first* superscript $`\mu `$ into $`\mu `$ in all components, while changing neither the second superscript nor the ordering of the fields.
Recall that the fields $`\varphi ^{(\mu ,\nu )}`$ are vectors with respect to the color group. Thus, the link multiplets $`\mathrm{\Phi }_a^i(n;\mu ),\mathrm{\Psi }_a^i(n;\mu )`$ have to be labeled by color ($`i`$) and flavor ($`a`$) indices. The latter are given by the second superscript in the definition of the multiplets: for instance, the flavor indices of the multiplets (20) take the successive values $`a=\mu +1,\mu 1`$, $`\mu +2,\mathrm{}`$ etc. Had we included multiple generations, the dimension of the coupling matrices would have been $`d^{}=n_b(2d2)`$).
Since there are $`2d`$ links around the site $`n`$ and for each link, two multiplets containing $`2d2`$ components, the full set of multiplets is of dimension $`8d(d1)`$. On the other hand, the number of independent flavor components at each site is $`2\times N_b=4d(d1)`$ (the factor $`2`$ corresponds to the chirality). Therefore, the site-link multiplets are not linearly independent; indeed, each field component $`\varphi _{L/R}^{(\pm \mu ,\pm \nu )}(n)`$ is contained in exactly two multiplets, one associated with the link $`(n+\widehat{\mu }/2)`$, the other with the link $`(n+\widehat{\nu }/2)`$.
In terms of these multiplets, the interacting part of the action (13) associated with the link $`(n+\widehat{\mu }/2)`$ can be written in a compact form as follows (repeated indices are summed over):
$$S_U(n+{\scriptscriptstyle \frac{\widehat{\mu }}{2}})=\overline{\mathrm{\Phi }}_a^i(n+\widehat{\mu };\mu )U^{ij}(n+{\scriptscriptstyle \frac{\widehat{\mu }}{2}})\mathrm{\Phi }_a^j(n;\mu )+\overline{\mathrm{\Psi }}_a^i(n;\mu )U^{ij}(n+{\scriptscriptstyle \frac{\widehat{\mu }}{2}})\mathrm{\Psi }_a^j(n+\widehat{\mu };\mu ).$$
(21)
Now that we isolated the part of the action associated with the matrix $`U(n+{\scriptscriptstyle \frac{\widehat{\mu }}{2}})`$, we can apply the bosonic $`U(N_c)`$ color-flavor transformation on this action, that is replace the integration over $`U(n+{\scriptscriptstyle \frac{\widehat{\mu }}{2}})`$ by an integral over a complex matrix $`Z(n+{\scriptscriptstyle \frac{\widehat{\mu }}{2}})`$ of dimension $`d^{}`$ :
$$\begin{array}{c}_{\mathrm{U}(N_c)}dU(n+{\scriptscriptstyle \frac{\widehat{\mu }}{2}})\mathrm{exp}[S_U(n+{\scriptscriptstyle \frac{\widehat{\mu }}{2}})]=_{D_d^{}}d\mu (Z,Z^{})\mathrm{Det}(1ZZ^{})^{N_c}\times \hfill \\ \hfill \times \mathrm{exp}\left[\overline{\mathrm{\Phi }}_a^i(n+\widehat{\mu };\mu )Z_{ab}(n+{\scriptscriptstyle \frac{\widehat{\mu }}{2}})\mathrm{\Psi }_b^i(n+\widehat{\mu };\mu )+\overline{\mathrm{\Psi }}_b^i(n;\mu )Z_{}^{}{}_{ba}{}^{}(n+{\scriptscriptstyle \frac{\widehat{\mu }}{2}})\mathrm{\Phi }_a^i(n;\mu )\right].\end{array}$$
(22)
$`D_d^{}`$ denotes the set of complex matrices $`Z`$ of dimension $`d^{}`$ such that the Hermitian matrix $`1ZZ^{}`$ is positive definite. This set is in one-to-one correspondence with the non-compact symmetric space $`\mathrm{U}(d^{},d^{})/\mathrm{U}(d^{})\times \mathrm{U}(d^{})`$ , and the measure $`d\mu (Z,Z^{})`$ is the (suitablly normalized) invariant measure on this symmetric space:
$$d\mu (Z,Z^{})=const\times \mathrm{Det}(1ZZ^{})^{2d^{}}\underset{a,b=1}{\overset{d^{}}{}}dZ_{ab}d\overline{Z}_{ab}.$$
The identity (22) makes sense iff
$$N_c2d^{}=4(d1),$$
(23)
otherwise the integral over $`Z`$ does not converge.
In the color-flavor transformed action, the auxiliary fields are coupled *ultralocally* via the $`Z`$ matrices through their flavor indices. The indices of the matrix $`Z(n+{\scriptscriptstyle \frac{\widehat{\mu }}{2}})`$ are associated with the plaquettes adjacent to the link $`(n+\widehat{\mu }/2)`$, so that each entry of that matrix describes a correlation between these plaquettes. This is to be put in contrast with the original action (13), which described a parallel transport of the bosonic field along the links. In other words, if the $`U`$-field is responsible for the transport along the link, the Z-field is responsible for the correlations between different plaquettes. We shall see that the latter correlations become suppressed in the large mass limit.
Grouping all auxiliary fields at the site $`n`$ we get the interaction part of the local effective action
$$S_Z[n]=\underset{\mu =1}{\overset{d}{}}\left[\overline{\mathrm{\Psi }}_b^i(n;\mu )\overline{Z}_{ab}(n+{\scriptscriptstyle \frac{\widehat{\mu }}{2}})\mathrm{\Phi }_a^i(n;\mu )+\overline{\mathrm{\Phi }}_a^i(n;\mu )Z_{ab}(n{\scriptscriptstyle \frac{\widehat{\mu }}{2}})\mathrm{\Psi }_b^i(n;\mu )\right],$$
(24)
which is diagonal in the color indices. Since the mass term is diagonal with respect to both the color and flavor indices, the color degrees of freedom are decoupled in the transformed action, so that the partition function can be factorized into $`N_c`$ identical integrals, each corresponding to a given color. Still, the coupling between the $`4d(d1)`$ flavor components at each site is not completely obvious, so we first analyze the simpler case of $`d=2`$ before turning to the general case.
### 3.1 d=2 effective action
The simplest possible situation corresponds to the model placed on the 2-dimensional square lattice spanned by two orthogonal unit vectors $`\widehat{1}`$ and $`\widehat{2}`$. Let us consider the four links having the lattice site $`n`$ in common (see Fig 3). The space of auxiliary fields at $`n`$ is of dimension $`8`$ and, according to (24,20), the site-link multiplets (here, doublets) read
$$\begin{array}{cc}\hfill \mathrm{\Psi }(n;1)& =\left(\genfrac{}{}{0pt}{}{\varphi _R^{(1,2)}(n)}{\varphi _L^{(1,2)}(n)}\right);\mathrm{\Psi }(n;1)=\left(\genfrac{}{}{0pt}{}{\varphi _R^{(1,2)}(n)}{\varphi _L^{(1,2)}(n)}\right)\hfill \\ \hfill \mathrm{\Psi }(n;2)& =\left(\genfrac{}{}{0pt}{}{\varphi _R^{(2,1)}(n)}{\varphi _L^{(2,1)}(n)}\right);\mathrm{\Psi }(n;2)=\left(\genfrac{}{}{0pt}{}{\varphi _R^{(2,1)}(n)}{\varphi _L^{(2,1)}(n)}\right).\hfill \end{array}$$
(25)
The $`4`$ chirally conjugated doublets $`\mathrm{\Phi }(n;\pm \widehat{\alpha })`$ are obtained by exchanging $`LR`$. These doublets are coupled through the $`2\times 2`$ matrices $`Z^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}}),Z^{}(n+{\scriptscriptstyle \frac{\widehat{2}}{2}}),Z(n{\scriptscriptstyle \frac{\widehat{1}}{2}})`$ and $`Z(n{\scriptscriptstyle \frac{\widehat{2}}{2}})`$ carried by the four links adjacent to the site $`n`$. To give an example, the matrix $`Z^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})`$ has the following index structure:
$$Z^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})=\left(\begin{array}{cc}Z_{2,2}^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})& Z_{2,2}^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})\\ Z_{2,2}^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})& Z_{2,2}^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})\end{array}\right)$$
(26)
Together with the link carrying the matrix, the lower pair of indices represent the plaquettes associated with the field components coupled by the matrix element: the diagonal element $`Z_{2,2}^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})`$ couples different fields associated with the same plaquette $`(n;1,2)`$, while the nondiagonal element $`Z_{}^{}{}_{2,2}{}^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})`$ couples fields associated with the two plaquettes $`(n;1,2)`$ and $`(n;1,2)`$.
We want to write an effective action uniquely in terms of the $`Z`$ fields, by integrating over the bosonic fields. For this aim, we need to describe the coupling between each pair or flavors in the action (24). As we already mentioned, the site-link multiplets (25) are not independent of each other, so we now group the auxiliary fields at the lattice site $`n`$ into chirally conjugated *site quadruplets*:
$$\mathrm{\Phi }(n)\stackrel{\mathrm{def}}{=}\left(\begin{array}{c}\varphi _R^{(1,2)}(n)\\ \varphi _L^{(1,2)}(n)\\ \varphi _L^{(1,2)}(n)\\ \varphi _R^{(1,2)}(n)\end{array}\right);\mathrm{\Psi }(n)\stackrel{\mathrm{def}}{=}\left(\begin{array}{c}\varphi _L^{(1,2)}(n)\\ \varphi _R^{(1,2)}(n)\\ \varphi _R^{(1,2)}(n)\\ \varphi _L^{(1,2)}(n)\end{array}\right).$$
(27)
The union of these two quadruplets contain each bosonic component once. The color-flavor transformed action (24) can be written in terms of these quadruplets via two complex $`4\times 4`$ matrices in the flavor space, $`V(n)`$ and $`W(n)`$, which contain the components of the $`Z`$-fields:
$$S_Z[n]=\mathrm{\Phi }^{}(n)V(n)\mathrm{\Psi }(n)+\mathrm{\Psi }^{}(n)W(n)\mathrm{\Phi }(n).$$
(28)
The matrices $`V(n)`$ and $`W(n)`$ can be compactly written
$$V(n)\stackrel{\mathrm{def}}{=}\left(\begin{array}{cc}Z^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})& 0\\ 0& Z(n{\scriptscriptstyle \frac{\widehat{1}}{2}})\end{array}\right);W(n)\stackrel{\mathrm{def}}{=}\tau _{(1,4)}\left(\begin{array}{cc}Z(n{\scriptscriptstyle \frac{\widehat{2}}{2}})& 0\\ 0& Z^{}(n+{\scriptscriptstyle \frac{\widehat{2}}{2}})\end{array}\right)\tau _{(1,4)}$$
(29)
where the permutation matrix $`\tau _{(1,4)}`$ interchanges the first and fourth indices. In other words, we make use of the basis elements of the simple matrix algebra of order 2
$$T_1e_{11}=\left(\begin{array}{cc}1& 0\\ 0& 0\end{array}\right)T_2e_{22}=\left(\begin{array}{cc}0& 0\\ 0& 1\end{array}\right),T_3e_{21}=\left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right),T_4e_{12}=\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right)$$
(30)
which satisfy
$$e_{ij}e_{kl}=\delta _{jk}e_{il},i,i,k,l=1,2$$
Then the matrices $`V(n)`$ and $`W(n)`$ can be expressed via the direct tensor products of first two of these generators and $`Z`$-fields for the ‘$`\widehat{1}`$’ and ‘$`\widehat{2}`$’ directed links respectively:
$$V\stackrel{\mathrm{def}}{=}T_1Z^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})+T_2(n{\scriptscriptstyle \frac{\widehat{1}}{2}});W\stackrel{\mathrm{def}}{=}\tau _{(1,4)}[T_1Z(n{\scriptscriptstyle \frac{\widehat{2}}{2}})+T_2Z^{}(n+{\scriptscriptstyle \frac{\widehat{2}}{2}})]\tau _{(1,4)}^1$$
(31)
where the matrix
$$\tau _{(1,4)}=\tau _{(1,4)}^1=T_1T_2+T_2T_1+T_3T_3+T_4T_4\left(\begin{array}{cccc}0& 0& 0& 1\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ 1& 0& 0& 0\end{array}\right)$$
(32)
permutes the first and the fourth components of the site quadruplets (27).
The integral over auxiliary fields at the site $`n`$ (including one color component) reads
$$\begin{array}{cc}\hfill 𝒵[n]=& 𝑑\mathrm{\Psi }^{}(n)𝑑\mathrm{\Psi }(n)𝑑\mathrm{\Phi }^{}(n)𝑑\mathrm{\Phi }(n)\mathrm{exp}\left[m_b(\mathrm{\Psi }^{}\mathrm{\Psi }+\mathrm{\Phi }^{}\mathrm{\Phi })+\mathrm{\Phi }^{}V\mathrm{\Psi }+\mathrm{\Psi }^{}W\mathrm{\Phi }\right]\hfill \\ & Det\left(\begin{array}{cc}m_b& V\\ W& m_b\end{array}\right)^1=Det(m_b^2VM)^1\hfill \\ & =\mathrm{exp}\left[\mathrm{Tr}\mathrm{ln}(1m_b^2VW)\right]\mathrm{exp}\left[m_b^2\mathrm{Tr}\left(VW\right)\right].\hfill \end{array}$$
(33)
In the last line we expanded the logarithm to first order in $`1/m_b`$.
The trace of the product $`VW`$ can be easily computed:
$$\begin{array}{c}\mathrm{Tr}\left(V(n)W(n)\right)=Z_{2,2}^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})Z_{1,1}^{}(n+{\scriptscriptstyle \frac{\widehat{2}}{2}})+Z_{2,2}^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})Z_{1,1}(n{\scriptscriptstyle \frac{\widehat{2}}{2}})\hfill \\ \hfill +Z_{2,2}(n{\scriptscriptstyle \frac{\widehat{1}}{2}})Z_{1,1}^{}(n+{\scriptscriptstyle \frac{\widehat{2}}{2}})+Z_{2,2}(n{\scriptscriptstyle \frac{\widehat{1}}{2}})Z_{1,1}(n{\scriptscriptstyle \frac{\widehat{2}}{2}}).\end{array}$$
(34)
Notice that only the diagonal elements of the $`Z`$ matrices appear in this leading-order term, which represent couplings between auxiliary fields carried by the same plaquette. In each term of the sum (34), the two matrix elements are carried by different links, but they correspond to fields related to the same plaquette, precisely the plaquette which shares these two links. One can represent the correlations embodied in (34) by links of *dual lattice* joining the middles of the two coupled links of the original lattice (see Fig. 4).
Clearly, the two-dimensional model is self-dual because both the original and dual lattices are square ones.
To summarize, to leading order in $`1/m_b`$ the full partition function is given by:
$$𝒵=\left\{\underset{n}{}\underset{\alpha =1,2}{}d\mu (Z,Z^{}(n+{\scriptscriptstyle \frac{\widehat{\alpha }}{2}}))\right\}\mathrm{exp}(N_cS[Z]),$$
(35)
with the effective action depending on the “flavor” matrices $`Z`$:
$$S[Z]=\underset{n}{}\left[m_b^2\mathrm{Tr}\left(V(n)W(n)\right)+\underset{\alpha =1,2}{}\mathrm{Tr}\mathrm{ln}\left(1Z(n+{\scriptscriptstyle \frac{\widehat{\alpha }}{2}})Z^{}(n+{\scriptscriptstyle \frac{\widehat{\alpha }}{2}})\right)\right].$$
(36)
### 3.2 Saddle point equations and continuum limit of the $`d=2`$ effective theory
The appearance of the factor $`N_c`$ in the exponent of the partition function (35) naturally suggests to imply the large-$`N_c`$ limit. Note that this limit is also needed to provide convergence of the bosonic color-flavor transformation, which is well-defined only if $`N_c2N_b=4d(d1)`$. Thus, in the color-flavor transformed model the large-$`N_c`$ limit becomes natural. Indeed, the structure of the action (36) suggests that in the large mass limit the functional integral over the $`Z,Z^{}`$ fields is sharply peaked about the matrices close to unity, that is we can make use of the scalar ansatz $`Z=z𝕀;Z^{}=z^{}𝕀`$, with $`z`$ and $`z^{}`$ two scalar functions to evaluate the partition function in the saddle point approximation.
After inserting this ansatz into (34) we obtain
$$\mathrm{Tr}\left(V(n)W(n)\right)=2[z^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})+z(n{\scriptscriptstyle \frac{\widehat{1}}{2}})][z^{}(n+{\scriptscriptstyle \frac{\widehat{2}}{2}})+z(n{\scriptscriptstyle \frac{\widehat{2}}{2}})]$$
(37)
and the action on the scalar ansatz becomes
$$S[z]=2\underset{n}{}\left[\underset{\alpha =1,2}{}\mathrm{ln}\left[1z(n+{\scriptscriptstyle \frac{\widehat{\alpha }}{2}})z^{}(n+{\scriptscriptstyle \frac{\widehat{\alpha }}{2}})\right]+m_{b}^{}{}_{}{}^{2}[z^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})+z(n{\scriptscriptstyle \frac{\widehat{1}}{2}})][z^{}(n+{\scriptscriptstyle \frac{\widehat{2}}{2}})+z(n{\scriptscriptstyle \frac{\widehat{2}}{2}})]\right].$$
(38)
In varying the action (38), the variables $`z`$ and $`z^{}`$ are to be considered as independent, which leads to two sets of saddle point equations
$$\begin{array}{cc}\hfill z^{}(n+{\scriptscriptstyle \frac{\widehat{2}}{2}})+z(n{\scriptscriptstyle \frac{\widehat{2}}{2}})& =\frac{m_b^2z(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})}{1zz^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})}=\frac{m_b^2z^{}(n{\scriptscriptstyle \frac{\widehat{1}}{2}})}{1zz^{}(n{\scriptscriptstyle \frac{\widehat{1}}{2}})};\hfill \\ \hfill z^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})+z(n{\scriptscriptstyle \frac{\widehat{1}}{2}})& =\frac{m_b^2z(n+{\scriptscriptstyle \frac{\widehat{2}}{2}})}{1zz^{}(n+{\scriptscriptstyle \frac{\widehat{2}}{2}})}=\frac{m_b^2z^{}(n{\scriptscriptstyle \frac{\widehat{2}}{2}})}{1zz^{}(n{\scriptscriptstyle \frac{\widehat{2}}{2}})},\hfill \end{array}$$
which can be resolved if $`z`$ is just a spacetime–independent real number:
$$z=z^{}=\sqrt{1\frac{m_b^2}{2}}$$
(39)
There is also a trivial solution $`z=z^{}=0`$.
Thus, the auxiliary scalar field cannot be set too heavy since the effective action on the saddle point configuration
$$S_{saddle}^{d=2}=2\left(\mathrm{ln}\frac{m_b^2}{2}1+\frac{2}{m_b^2}\right)$$
becomes imaginary as $`m_b^2>2`$ in units of lattice spacing $`a`$. However, even if we restrict the mass of inducing fields to the interval $`1<m_b<2`$, the corresponding induced coupling $`\beta `$ remains small enough to justify the large mass expansion.
Let us consider a continuum limit of the model (36). We again use the scalar ansatz for the Z-fields to write the action in the form (38).
Following , we can redefine the variables $`z(n+{\scriptscriptstyle \frac{\widehat{1}}{2}}),z(n+{\scriptscriptstyle \frac{\widehat{2}}{2}})`$, which appears in the *middle* of the links $`\widehat{1},\widehat{2}`$, as differences of $`\omega (n)`$ taken at the neighboring sites:
$$\begin{array}{cc}\hfill z(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})& =\omega (n+1,2)\omega (1,2)a_x\omega (n)+\frac{a^2}{2}_x^2\omega (n)+\mathrm{};\hfill \\ \hfill z(n{\scriptscriptstyle \frac{\widehat{1}}{2}})& =\omega (1,2)\omega (n1,2)a_x\omega (n)\frac{a^2}{2}_x^2\omega (n)+\mathrm{};\hfill \end{array}$$
etc. Here the coordinate $`x`$ is taken along positive direction of the link $`\widehat{1}`$ and the coordinate $`y`$ is taken alon positive direction of the link $`\widehat{2}`$.
The Jacobian of this transformation is given by the determinant of $`2\times 2`$ matrix composed of the second derivatives: $`J(\omega )=\mathrm{Det}(_x_y\omega )`$. In terms of these new variables one easily obtains
$$\begin{array}{cc}\hfill \underset{\alpha =\widehat{1},\widehat{2}}{}\mathrm{ln}\left[1z(n+{\scriptscriptstyle \frac{\widehat{\alpha }}{2}})z^{}(n+{\scriptscriptstyle \frac{\widehat{\alpha }}{2}})\right]\underset{\alpha =x,y}{}\mathrm{ln}(1+a^2_\alpha \omega _\alpha \omega ^{})& a^2\underset{\alpha =x,y}{}(_\alpha \omega )(_\alpha \omega ^{});\hfill \\ \hfill [z^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})+z(n{\scriptscriptstyle \frac{\widehat{1}}{2}})][z^{}(n+{\scriptscriptstyle \frac{\widehat{2}}{2}})+z(n{\scriptscriptstyle \frac{\widehat{2}}{2}})]& 2a^2\mathrm{Re}_x\omega _y\omega .\hfill \end{array}$$
Recall that in the $`d=2`$ model the correlation length $`\xi =m_b^1`$ remains finite in the continuum limit. Then the effective continuum coupling $`\beta =m_b^2`$ also is finite and in the continuum limit the effective two-dimensional action becomes
$$\begin{array}{cc}\hfill S& =a^2\underset{n}{}\left(\underset{\alpha =x,y}{}(_\alpha \omega )(_\alpha \omega ^{})\right)+2m_b^2\mathrm{Re}_x\omega _y\omega \hfill \\ & d^2x\left[(_x\omega )(_x\omega ^{})+(_y\omega )(_y\omega ^{})+2\beta \mathrm{Re}_x\omega _y\omega \right]\hfill \end{array}$$
(40)
where we restore conventional coordinate notations $`\mu x`$ and $`\nu y`$ in the continuum form of the effective action. Expression (40) represents a dual version of d=2 inducing theory in the large-$`N_c`$ strong coupling limit. Recall that the $`Z`$ field appears after the color-flavor transformation and therefore the action (40) by definition does not contain color degrees of freedom. One may consider it as the action of the $`\sigma `$-model perturbed by the non-relativistic interaction term with a coupling constant $`\beta `$.
Furthermore, we can use the method of Ref. and identify the functions $`\omega `$ with the external coordinates of some manifold. The metric tensor of the manifold is $`g_{ij}=_i\omega _j\omega ^{}`$ and the determinant of it is
$$\mathrm{Det}g_{ij}g=\frac{1}{2}\epsilon ^{ij}\epsilon ^{kl}(_i\omega _k\omega ^{})(_j\omega _l\omega ^{}).$$
Then the kinetic term of the continuum action (40) can be represented in terms of the metric tensor as $`(_i\omega )(_i\omega ^{})=g_{ii}`$ in correspondence with Ref. . However, there is no straihgtforward generalization of the above considered continuum limit of $`d=2`$ model for a higher dimensional case.
### 3.3 $`d=3`$ effective action
The consideration of the model on 3-dimensional lattice spanned by three vectors $`\widehat{\mu },\widehat{\nu },\widehat{\rho }`$ becomes more complicated since the space of auxiliary fields is now of dimension of 24. First, we have to fix the orientation of the plaquettes in 3 planes $`(\mu \nu ),(\mu \rho )`$ and $`(\nu \rho )`$ according to the convention (8). We suppose that $`\mu 1`$, $`\nu 2`$ and $`\rho 3`$, thus the left chiral components run on these plaquettes as in Fig. 5.
Let us analyze the structure of the action (24) in that particular case. According to the definition of the link multiplets (20), there are 6 scalar quadruplets whose explicit form is
$$\begin{array}{cc}\hfill \mathrm{\Psi }(n;1)& \stackrel{\mathrm{def}}{=}\left(\begin{array}{c}\varphi _R^{(1,2)}(n)\\ \varphi _L^{(1,2)}(n)\\ \varphi _R^{(1,3)}(n)\\ \varphi _L^{(1,3)}(n)\end{array}\right);\mathrm{\Psi }(n;2)\stackrel{\mathrm{def}}{=}\left(\begin{array}{c}\varphi _R^{(2,3)}(n)\\ \varphi _L^{(2,3)}(n)\\ \varphi _R^{(2,1)}(n)\\ \varphi _L^{(2,1)}(n)\end{array}\right);\mathrm{\Psi }(n;3)\stackrel{\mathrm{def}}{=}\left(\begin{array}{c}\varphi _R^{(3,1)}(n)\\ \varphi _L^{(3,1)}(n)\\ \varphi _R^{(3,2)}(n)\\ \varphi _L^{(3,2)}(n)\end{array}\right);\hfill \\ \hfill \mathrm{\Phi }(n;1)& \stackrel{\mathrm{def}}{=}\left(\begin{array}{c}\varphi _L^{(1,2)}(n)\\ \varphi _R^{(1,2)}(n)\\ \varphi _L^{(1,3)}(n)\\ \varphi _R^{(1,3)}(n)\end{array}\right);\mathrm{\Phi }(n;2)\stackrel{\mathrm{def}}{=}\left(\begin{array}{c}\varphi _L^{(2,3)}(n)\\ \varphi _R^{(2,3)}(n)\\ \varphi _L^{(2,1)}(n)\\ \varphi _R^{(2,1)}(n)\end{array}\right);\mathrm{\Phi }(n;3)\stackrel{\mathrm{def}}{=}\left(\begin{array}{c}\varphi _L^{(3,1)}(n)\\ \varphi _R^{(3,1)}(n)\\ \varphi _L^{(3,2)}(n)\\ \varphi _R^{(3,2)}(n)\end{array}\right),\hfill \end{array}$$
These multiplets are coupled with 6 matrix $`Z`$-fields of dimension $`4\times 4`$ defined on the 6 directed links $`\widehat{1}`$, $`\widehat{2}`$ and $`\widehat{3}`$ which are adjacent to the point $`n`$:
$$S_Z[\mathrm{\Psi }(n),\mathrm{\Phi }(n)]=\underset{\alpha =1}{\overset{3}{}}\left[\overline{\mathrm{\Psi }}(n;\alpha )Z^{}(n+{\scriptscriptstyle \frac{\widehat{\alpha }}{2}})\mathrm{\Phi }(n;\alpha )+\overline{\mathrm{\Phi }}(n;\alpha )Z(n{\scriptscriptstyle \frac{\widehat{\alpha }}{2}})\mathrm{\Psi }(n;\alpha )\right]$$
(41)
The explicit form of the matrix $`Z^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})`$ for example is (cf. its $`2d`$ counterpart (26))
$$Z^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})=\left(\begin{array}{cccc}Z_{2,2}^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})& Z_{2,2}^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})& Z_{2,3}^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})& Z_{2,3}^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})\\ Z_{2,2}^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})& Z_{2,2}^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})& Z_{2,3}^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})& Z_{2,3}^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})\\ Z_{3,2}^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})& Z_{3,2}^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})& Z_{3,3}^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})& Z_{3,3}^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})\\ Z_{3,2}^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})& Z_{3,2}^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})& Z_{3,3}^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})& Z_{3,3}^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})& \end{array}\right),\mathrm{etc}$$
(42)
To simplify the evaluation of the functional determinant of the color-flavor transformed model in $`d=3`$ we have to bring the matrix, which couples all these multiplets and contains different components of the $`Z`$ fields, to the block-diagonal form. Thus we again regroup the components of the scalar link multiplets (3.3) into the site quadruplets:
$$\begin{array}{cc}\hfill \mathrm{\Phi }^{(12)}& \stackrel{\mathrm{def}}{=}\left(\begin{array}{c}\varphi _R^{(1,2)}\\ \varphi _L^{(1,2)}\\ \varphi _L^{(1,2)}\\ \varphi _R^{(1,2)}\end{array}\right);\mathrm{\Phi }^{(13)}\stackrel{\mathrm{def}}{=}\left(\begin{array}{c}\varphi _R^{(1,3)}\\ \varphi _L^{(1,3)}\\ \varphi _L^{(1,3)}\\ \varphi _R^{(1,3)}\end{array}\right);\mathrm{\Phi }^{(23)}\stackrel{\mathrm{def}}{=}\left(\begin{array}{c}\varphi _R^{(2,3)}\\ \varphi _L^{(2,3)}\\ \varphi _L^{(2,3)}\\ \varphi _R^{(2,3)}\end{array}\right)\hfill \\ \hfill \mathrm{\Psi }^{(12)}& \stackrel{\mathrm{def}}{=}\left(\begin{array}{c}\varphi _L^{(1,2)}\\ \varphi _R^{(1,2)}\\ \varphi _R^{(1,2)}\\ \varphi _L^{(1,2)}\end{array}\right);\mathrm{\Psi }^{(13)}\stackrel{\mathrm{def}}{=}\left(\begin{array}{c}\varphi _L^{(1,3)}\\ \varphi _R^{(1,3)}\\ \varphi _R^{(1,3)}\\ \varphi _L^{(1,3)}\end{array}\right);\mathrm{\Psi }^{(23)}\stackrel{\mathrm{def}}{=}\left(\begin{array}{c}\varphi _L^{(2,3)}\\ \varphi _R^{(2,3)}\\ \varphi _R^{(2,3)}\\ \varphi _L^{(2,3)}\end{array}\right)\hfill \end{array}$$
(43)
In terms of these variables the local actions (41) can be rewritten as (cf. Eq. (28))
$$\begin{array}{cc}\hfill S_Z[n]& =\overline{\mathrm{\Phi }}^{(12)}V_{(1)}^{(2,2)}\mathrm{\Psi }^{(12)}+\overline{\mathrm{\Phi }}^{(13)}V_{(1)}^{(3,3)}\mathrm{\Psi }^{(13)}+\overline{\mathrm{\Phi }}^{(12)}V_{(1)}^{(2,3)}\mathrm{\Psi }^{(13)}\hfill \\ & +\overline{\mathrm{\Phi }}^{(13)}V_{(1)}^{(3,2)}\mathrm{\Psi }^{(12)}+\overline{\mathrm{\Phi }}^{(23)}V_{(2)}^{(3,3)}\mathrm{\Psi }^{(23)}+\overline{\mathrm{\Psi }}^{(12)}W_{(2)}^{(1,1)}\mathrm{\Phi }^{(12)}\hfill \\ & +\overline{\mathrm{\Psi }}^{(13)}W_{(3)}^{(1,1)}\mathrm{\Phi }^{(13)}+\overline{\mathrm{\Psi }}^{(13)}W_{(3)}^{(1,2)}\mathrm{\Phi }^{(23)}+\overline{\mathrm{\Psi }}^{(23)}W_{(3)}^{(2,1)}\mathrm{\Phi }^{(13)}\hfill \\ & +\overline{\mathrm{\Psi }}^{(23)}W_{(3)}^{(2,2)}\mathrm{\Phi }^{(23)}+\overline{\mathrm{\Phi }}^{(23)}R_{(2)}^{(3,1)}\mathrm{\Phi }^{(12)}+\overline{\mathrm{\Psi }}^{(12)}S_{(2)}^{(1,3)}\mathrm{\Psi }^{(23)}\hfill \end{array}$$
where the components of the $`Z`$ fields are now regrouped into new $`4\times 4`$ matrices which are coupled with the site quadruplets (43) in the action (3.3). The lower index of these matrix fields is associated with the link which carries the components of the $`Z`$ field, while the upper pair represents the plaquette indices of the site quadruplets which cap the link.
The matrices $`V_{(\alpha )}^{(\beta ,\gamma )},\alpha =1,2;\beta ,\gamma =2,3`$ are defined on the first two links and couple the multiplets $`\overline{\mathrm{\Phi }}`$ and $`\mathrm{\Psi }`$; one of these matrices is the above defined d=2 couplig matrix $`V(n)V_{(1)}^{(2,2)}`$ which appears in (33). The matrices $`W_{(\alpha )}^{(\beta ,\gamma )},\alpha =2,3;\beta ,\gamma =1,2`$ are defined on the second and third links and couple the multiplets $`\overline{\mathrm{\Psi }}`$ and $`\mathrm{\Phi }`$, one of these matrices also appears in (33). Unlike the $`d=2`$ model there are two more matrices $`R_{(2)}^{(3,1)}`$ and $`S_{(2)}^{(1,3)}`$ which live on the second links and couple the fields $`\overline{\mathrm{\Phi }},\mathrm{\Phi }`$ and $`\overline{\mathrm{\Psi }},\mathrm{\Psi }`$, respectively.
These matrices can be written concisely by making use of the above defined generators $`T_i`$ (30). Let us supplement the permutation matrix $`\tau _2`$ of eq. (32) by two other matrices
$$\begin{array}{cc}\hfill \tau _1& =T_3𝕀++T_4𝕀\left(\begin{array}{cccc}0& 0& 1& 0\\ 0& 0& 0& 1\\ 1& 0& 0& 0\\ 0& 1& 0& 0\end{array}\right);\hfill \\ \hfill \tau _3& =T_1T_4+T_4T_2+T_2T_3+T_3T_1\left(\begin{array}{cccc}0& 1& 0& 0\\ 0& 0& 0& 1\\ 1& 0& 0& 0\\ 0& 0& 1& 0\end{array}\right),\hfill \end{array}$$
which permute all four indices of the site quadruplets simultaneously, $`(1,2,3,4)(3,4,1,2)`$ and $`(1,2,3,4)(3,1,4,2)`$ respectively.
One can consider the basis elements of the matrix algebra (30) as generators whose tensor products with a complex $`2\times 2`$ matrix define different embeddings into a matrix of dimension $`4\times 4`$. Then each of the 6 matrices of the $`Z`$ field can be written as an expansion in the basis of the generators $`T_i`$:
$$Z=\underset{i=1}{\overset{4}{}}Z_iT_i$$
where $`2\times 2`$ blocks in components are given by
$$Z_1=\left(\begin{array}{cc}Z_{11}& Z_{12}\\ Z_{21}& Z_{22}\end{array}\right);Z_2=\left(\begin{array}{cc}Z_{33}& Z_{34}\\ Z_{43}& Z_{44}\end{array}\right);Z_3=\left(\begin{array}{cc}Z_{31}& Z_{32}\\ Z_{41}& Z_{42}\end{array}\right);Z_4=\left(\begin{array}{cc}Z_{13}& Z_{14}\\ Z_{23}& Z_{24}\end{array}\right)$$
(44)
With all these shortnotes at hand, we can write
$$\begin{array}{cc}\hfill V_{(1)}^{(2,2)}& =Z_1^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})+\tau _1Z_1(n{\scriptscriptstyle \frac{\widehat{1}}{2}})\tau _1^1;V_{(1)}^{(3,3)}=\tau _1Z_2^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})\tau _1^1+Z_2(n{\scriptscriptstyle \frac{\widehat{1}}{2}});\hfill \\ \hfill V_{(1)}^{(2,3)}& =Z_4^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})\tau _1^1+\tau _1Z_4(n{\scriptscriptstyle \frac{\widehat{1}}{2}});V_{(1)}^{(3,2)}=\tau _1Z_3^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})+Z_3(n{\scriptscriptstyle \frac{\widehat{1}}{2}})\tau _1^1;\hfill \\ \hfill V_{(2)}^{(3,3)}& =Z_1^{}(n+{\scriptscriptstyle \frac{\widehat{2}}{2}})+\tau _1Z_1(n{\scriptscriptstyle \frac{\widehat{2}}{2}})\tau _1^1;W_{(2)}^{(1,1)}=\tau _2Z_2^{}(n+{\scriptscriptstyle \frac{\widehat{2}}{2}})\tau _2^1+\tau _3Z_2(n{\scriptscriptstyle \frac{\widehat{2}}{2}})\tau _3^1;\hfill \\ \hfill W_{(3)}^{(1,1)}& =\tau _3Z_1^{}(n+{\scriptscriptstyle \frac{\widehat{3}}{2}})\tau _3^1+\tau _2Z_1(n{\scriptscriptstyle \frac{\widehat{3}}{2}})\tau _2^1;W_{(3)}^{(2,2)}=\tau _2Z_2^{}(n+{\scriptscriptstyle \frac{\widehat{3}}{2}})\tau _2^1+\tau _3Z_2(n{\scriptscriptstyle \frac{\widehat{3}}{2}})\tau _3^1;\hfill \\ \hfill W_{(3)}^{(1,2)}& =\tau _3Z_4^{}(n+{\scriptscriptstyle \frac{\widehat{3}}{2}})\tau _2^1+\tau _2Z_4(n{\scriptscriptstyle \frac{\widehat{3}}{2}})\tau _3^1;W_{(3)}^{(2,1)}=\tau _2Z_3^{}(n+{\scriptscriptstyle \frac{\widehat{3}}{2}})\tau _3^1+\tau _3Z_3(n{\scriptscriptstyle \frac{\widehat{3}}{2}})\tau _2^1;\hfill \\ \hfill S_{(2)}^{(1,3)}& =\tau _2Z_3^{}(n+{\scriptscriptstyle \frac{\widehat{2}}{2}})+\tau _3Z_3(n{\scriptscriptstyle \frac{\widehat{2}}{2}})\tau _1^1;R_{(2)}^{(3,1)}=Z_4^{}(n+{\scriptscriptstyle \frac{\widehat{2}}{2}})\tau _2^1+\tau _1Z_4(n{\scriptscriptstyle \frac{\widehat{2}}{2}})\tau _3^1\hfill \end{array}$$
To complete our calculations and carry out the integration over the scalar fields, we compose the site multiplets (43) into 12-plets
$$\mathrm{\Phi }\stackrel{\mathrm{def}}{=}\left(\begin{array}{c}\mathrm{\Phi }^{(12)}\\ \mathrm{\Phi }^{(13)}\\ \mathrm{\Phi }^{(23)}\end{array}\right);\mathrm{\Psi }\stackrel{\mathrm{def}}{=}\left(\begin{array}{c}\mathrm{\Psi }^{(12)}\\ \mathrm{\Psi }^{(13)}\\ \mathrm{\Psi }^{(23)}\end{array}\right)$$
which become coupled with the total coupling matrices. These matrices have the following block structure
$$V\stackrel{\mathrm{def}}{=}\left(\begin{array}{ccc}V_{(1)}^{(2,2)}& V_{(1)}^{(3,2)}& 0\\ V_{(1)}^{(2,3)}& V_{(1)}^{(3,3)}& 0\\ 0& 0& V_{(2)}^{(3,3)}\end{array}\right);W\stackrel{\mathrm{def}}{=}\left(\begin{array}{ccc}W_{(2)}^{(1,1)}& 0& 0\\ 0& W_{(3)}^{(1,1)}& W_{(3)}^{(1,2)}\\ 0& W_{(3)}^{(2,1)}& W_{(3)}^{(2,2)}\end{array}\right).$$
(45)
Clearly, the mass term of the auxiliary fields needs to be taken into account in the final expression for the action. Then the matrices $`S_{(2)}^{(1,3)}`$ and $`R_{(2)}^{(3,1)}`$ appear as off-diagonal blocks of the total mass matrices
$$M_1\stackrel{\mathrm{def}}{=}\left(\begin{array}{ccc}m_b& 0& 0\\ 0& m_b& 0\\ R_{(2)}^{(3,1)}& 0& m_b\end{array}\right);M_2\stackrel{\mathrm{def}}{=}\left(\begin{array}{ccc}m_b& 0& S_{(2)}^{(1,3)}\\ 0& m_b& 0\\ 0& 0& m_b\end{array}\right)$$
(46)
respectively. Note that $`DetM_1=DetM_2=m_b^3`$.
Thus, the complete action can be written in the form
$$S[n]=\overline{\mathrm{\Phi }}V\mathrm{\Psi }+\overline{\mathrm{\Psi }}W\mathrm{\Phi }+\overline{\mathrm{\Phi }}M_1\mathrm{\Phi }+\overline{\mathrm{\Psi }}M_2\mathrm{\Psi }$$
and the functional integration over the auxiliary fields yields the determinant of the full coupling matrix of dimension 24:
$$Det^{N_c}\left(\begin{array}{cc}M_1& V\\ W& M_2\end{array}\right)=Det^{N_c}(M_1M_2VW)Det^{N_c}(1M_2^1M_1^1VW)$$
(47)
Since
$$M_1^1=\left(\begin{array}{ccc}m_b^1& 0& 0\\ 0& m_b^1& 0\\ R_{(2)}^{(3,1)}m_b^2& 0& m_b^1\end{array}\right);M_2^1=\left(\begin{array}{ccc}m_b^1& 0& S_{(2)}^{(1,3)}m_b^2\\ 0& m_b^1& 0\\ 0& 0& m_b^1\end{array}\right),$$
the leading order of the large mass expansion yields
$$Det^{N_c}(1M_2^1M_1^1VW)\mathrm{exp}[N_c\mathrm{Tr}\mathrm{ln}(1\frac{1}{m_b^2}VW]\mathrm{exp}\left[\frac{N_c}{m_b^2}\mathrm{Tr}(VW)\right]$$
(48)
where the explicit form of the trace is rather cumbersome
$$\begin{array}{cc}\hfill \mathrm{Tr}\left(V(n)W(n)\right)& =Z_{2,2}^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})Z_{1,1}^{}(n+{\scriptscriptstyle \frac{\widehat{2}}{2}})+Z_{2,2}(n{\scriptscriptstyle \frac{\widehat{1}}{2}})Z_{1,1}^{}(n+{\scriptscriptstyle \frac{\widehat{2}}{2}})\hfill \\ & +Z_{2,2}^{}Z_{1,1}(n{\scriptscriptstyle \frac{\widehat{2}}{2}})+Z_{2,2}(n{\scriptscriptstyle \frac{\widehat{1}}{2}})Z_{1,1}(n{\scriptscriptstyle \frac{\widehat{2}}{2}})\hfill \\ & +Z_{1,1}^{}(n+{\scriptscriptstyle \frac{\widehat{3}}{2}})Z_{3,3}^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})+Z_{2,2}^{}(n+{\scriptscriptstyle \frac{\widehat{3}}{2}})Z_{3,3}(n{\scriptscriptstyle \frac{\widehat{1}}{2}})\hfill \\ & +Z_{3,3}^{}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})Z_{1,1}(n{\scriptscriptstyle \frac{\widehat{3}}{2}})+Z_{1,1}(n{\scriptscriptstyle \frac{\widehat{3}}{2}})Z_{3,3}(n{\scriptscriptstyle \frac{\widehat{1}}{2}})\hfill \\ & +Z_{2,2}^{}(n+{\scriptscriptstyle \frac{\widehat{3}}{2}})Z_{3,3}^{}(n+{\scriptscriptstyle \frac{\widehat{2}}{2}})+Z_{3,3}^{}(n+{\scriptscriptstyle \frac{\widehat{2}}{2}})Z_{2,2}(n{\scriptscriptstyle \frac{\widehat{3}}{2}})\hfill \\ & +Z^{}2,2(n+{\scriptscriptstyle \frac{\widehat{3}}{2}})Z_{3,3}(n{\scriptscriptstyle \frac{\widehat{2}}{2}})+Z_{2,2}(n{\scriptscriptstyle \frac{\widehat{3}}{2}})Z_{3,3}(n{\scriptscriptstyle \frac{\widehat{2}}{2}})\hfill \end{array}$$
(49)
Note, that as in the $`d=2`$ model (cf. Eq. (34), in the leading order of the expansion in $`1/m_b`$ only diagonal elements of the $`Z`$-matrix fields contribute to the functional determinant and the correlations between the plaquettes, which are described by (49) are factorized into the sum of 12 independent two-points correlations between the diagonal components of the $`Z`$ matrix field. Therefore, the cube of the original d=3 lattice on the dual lattice corresponds to the *tetradecahedra* whose edges corresponds to these correlations (see Fig. 6). The plaquette of the original lattice corresponds to a *tetrahedron*, six of those precisely fit the tetradecahedron of a dual cube. This is exactly the structure suggested in from a different viewpoint. However, the higher order corrections in $`1/m_b^2`$ make the geometry of the dual lattice much more complicated.
Finally, we arrive at the effective action of the 3d dual gluodynamics which has the same structure as (36):
$$S[Z]=\underset{n}{}\mathrm{Tr}\left[m_b^2V(n)W(n)+\underset{\alpha =1}{\overset{3}{}}\mathrm{ln}\left(1Z(n+{\scriptscriptstyle \frac{\widehat{\alpha }}{2}})Z^{}(n+{\scriptscriptstyle \frac{\widehat{\alpha }}{2}})\right)\right].$$
(50)
We will see that this results holds in any dimension.
Note that there is a relation between the space of the pairs of plaquette indices of $`d`$-dimensional space $`(\alpha ,\beta )`$ where $`\alpha ,\beta =1,2\mathrm{}d`$, and the Cartan–Weyl *simple root spaces* of two Lie $`𝔰𝔲(d)`$-algebras which are subalgebras of $`𝔤𝔩(d)`$. Indeed, one can set a correspondence between the ordered pairs $`(\alpha ,\beta )`$ with positive $`\alpha `$ and $`\beta `$, and the basis elements of the simple matrix algebra of $`d`$th order:
$$\begin{array}{cc}\hfill (\alpha ,\beta )& e_{\alpha \beta }\mathrm{if}\alpha <\beta ,\alpha ,\beta >0;\hfill \\ \hfill (\alpha ,\beta )& e_{\beta \alpha }\mathrm{if}\alpha <\beta ,\alpha ,\beta >0.\hfill \end{array}$$
Here the $`d\times d`$-matrices $`e_{\alpha \beta }`$ are d-dimensional generalization of the two-dimensional basis (30): the entry in the $`\alpha `$th row and $`\beta `$th column of $`e_{\alpha \beta }`$ is equal to 1 while all other entries are zero. With that ordering the matrices $`e_{\alpha \beta }`$ and $`e_{\beta \alpha }`$ can be identified as the raising and lowering generators $`E_{\pm \stackrel{}{\beta }_i}`$ of an $`𝔰𝔲(d)`$-algebra which is characterized by the set of the simple roots $`\pm \stackrel{}{\beta }_i`$ respectively. To cover all the space of the pairs of the plaquette indices we have to consider also the second set of the ordered pairs $`(\alpha ,\beta )`$ and $`(\alpha ,\beta )`$, which can be set into correspondence with a complimentary $`\overline{𝔰𝔲}(d)`$-algebra in the same way.
For example, in the d=3 model with the above ordering of the indices, the plaquettes $`(n;1,2)`$ and $`(n;1,2)`$ correspond to the matrices
$$(n;1,2)e_{12}=\left(\begin{array}{ccc}0& 1& 0\\ 0& 0& 0\\ 0& 0& 0\end{array}\right)E_{\stackrel{}{\beta }_1};(n;1,2)e_{21}=\left(\begin{array}{ccc}0& 0& 0\\ 1& 0& 0\\ 0& 0& 0\end{array}\right)E_{\stackrel{}{\beta }_1},$$
where $`\stackrel{}{\beta }_1`$ is one of the simple roots of the $`𝔰𝔲(3)`$-algebra. The second positive root $`\stackrel{}{\beta }_2`$ can be set into correspondence with the plaquette $`(n;2,3)`$ whereas the third positive composite root $`\stackrel{}{\beta }_3=\stackrel{}{\beta }_1+\stackrel{}{\beta }_2`$ corresponds to the plaquette $`(n;1,3)`$.
To identify the raising and lowering generators of the second $`\overline{𝔰𝔲}(3)`$-algebra we have to consider, for example, the pairs of the plaquettes $`(n;1,2)`$ and $`(n;1,2)`$ which can be mapped into the matrices $`\overline{e}_{12}`$ and $`\overline{e}_{21}`$ in the same way. This corresponds to the first simple root of the $`\overline{𝔰𝔲}(3)`$-algebra.
This construction can be used to integrate over the auxiliary fields in a general d-dimensional case .
### 3.4 Effective action in arbitrary dimension
Our aim in this section is the same as in the previous two, that is integrate the action (24) over the auxiliary fields at the site $`n`$, and compute the resulting effective action in the matrices $`Z`$, $`Z^{}`$ in arbitrary dimension. We will only consider the case of one generation of auxiliary fields, that is, $`n_b=1`$ . As we already pointed out, the difficulty comes from the fact that the same field $`\varphi _{L/R}^{(\pm \mu ,\pm \nu )}`$ is contained in two different site-links multiplets (20). In order to integrate over these fields, we first need to regroup them, that is write the action $`S_Z`$ as
$$S_Z[n]=\left[\varphi _{L/R}^{(\pm \mu ,\pm \nu )}\right]^{}(n)\left[\varphi _{L/R}^{(\pm \mu ,\pm \nu )}\right],$$
(51)
where the column vector $`\left[\varphi _{L/R}^{(\pm \mu ,\pm \nu )}\right]`$ contains the $`2N_b=4d(d1)`$ fields at the site $`n`$. The coupling matrix $`(n)`$ is therefore of dimension $`4d(d1)`$; it contains components of the matrices $`Z`$ and $`Z^{}`$ carried by the links touching $`n`$. Our task is to write down the matrix $``$ explicitly, using a judicious grouping of the field components. As was already the case in two dimensions, the matrix $``$ has many null entries, so that its determinant may be simplified.
We will group the auxiliary fields in *site quadruplets* associated with the planes in the $`d`$-dimensional lattice. Each plane, indexed by a couple $`(\mu \nu )`$ with $`1\mu <\nu d`$, contains $`4`$ plaquettes touching $`n`$. To each plane we associate two site quadruplets at $`n`$:
$$\mathrm{\Phi }^{(\mu \nu )}(n)\stackrel{\mathrm{def}}{=}\left(\begin{array}{c}\varphi _R^{(\mu ,\nu )}(n)\\ \varphi _L^{(\mu ,\nu )}(n)\\ \varphi _L^{(\mu ,\nu )}(n)\\ \varphi _R^{(\mu ,\nu )}(n)\end{array}\right);\mathrm{\Psi }^{(\mu \nu )}(n)\stackrel{\mathrm{def}}{=}\left(\begin{array}{c}\varphi _L^{(\mu ,\nu )}(n)\\ \varphi _R^{(\mu ,\nu )}(n)\\ \varphi _R^{(\mu ,\nu )}(n)\\ \varphi _L^{(\mu ,\nu )}(n)\end{array}\right).$$
(52)
These quadruplets generalize the ones defined in Eq. (27) to any plane in the $`d`$-dimensional lattice. The total number of these planes is $`\frac{d(d1)}{2}`$, so that the ‘concatenation’ of all the above site quadruplets yields the correct number of field components. To perform this concatenation, we need to *order* the different planes, that is, to order the couples $`(\mu \nu )`$.
These couples are in one-to-one correspondence with the positive roots of the Lie algebra $`𝔤𝔩(d)`$: each plane $`(\mu \nu )`$ can indeed be associated with the generator $`e_{\mu \nu }`$ of the algebra, satisfying the relations
$$[e_{\mu \nu },e_{\rho \eta }]=\delta _{\nu \rho }e_{\mu \eta }\delta _{\mu \eta }e_{\rho \nu }.$$
(53)
There is no canonical ordering of the positive generators (or the positive roots), on the other hand it seems natural to require that $`(\mu \nu )<(\rho \eta )`$ if $`\nu \rho `$; this condition is satisfied by the following ordering:
$$(12)<(13)<(14)<\mathrm{}<(1d)<(23)<(24)<\mathrm{}<(2d)<(34)<\mathrm{}<(d1d).$$
(54)
The site quadruplets will be ordered according to the above convention, starting with all quadruplets $`\mathrm{\Phi }^{(\mu \nu )}`$ and finishing with the quadruplets $`\mathrm{\Psi }^{(\mu \nu )}`$.
Now that we ordered the vector $`\left[\varphi _{L/R}^{(\pm \mu ,\pm \nu )}\right]`$, we need to compute the matrix $``$, and for this to derive which quadruplets $`\mathrm{\Phi }^{(\mu \nu )}`$ or $`\mathrm{\Psi }^{(\mu \nu )}`$ are coupled with which quadruplets $`\mathrm{\Phi }^{(\rho \eta )}`$ or $`\mathrm{\Psi }^{(\rho \eta )}`$, and through which matrices $`Z`$ or $`Z^{}`$. For this aim, we have to compare the components of, on one side, the quadruplets $`\mathrm{\Phi }^{(\mu \nu )}`$, $`\mathrm{\Psi }^{(\mu \nu )}`$; on the other side, the site-link multiplets $`\mathrm{\Phi }(n;\pm \alpha )`$, $`\mathrm{\Psi }(n;\pm \alpha )`$ which were used to write the action (24). For instance, the first component $`\varphi _R^{(\mu ,\nu )}`$ of $`\mathrm{\Phi }^{(\mu \nu )}`$ is contained in the site-link multiplet $`\mathrm{\Psi }(n;\mu )`$ (because $`\mu <\nu `$), so its complex conjugate is coupled to the matrix $`Z^{}(n+{\scriptscriptstyle \frac{\widehat{\mu }}{2}})`$ on the left; on the other hand, $`\varphi _R^{(\mu ,\nu )}`$ is also contained in the multiplet $`\mathrm{\Phi }(n;\nu )`$ (because $`\nu >\mu `$), so it is coupled to the matrix $`Z^{}(n+{\scriptscriptstyle \frac{\widehat{\nu }}{2}})`$ on the right. Below we schematically represent the couplings of the quadruplets with the $`Z`$ matrices by taking all components in the quadruplets into account:
$$\begin{array}{cc}\hfill \mathrm{\Phi }^{(\mu \nu )}Z^{}(n+{\scriptscriptstyle \frac{\widehat{\mu }}{2}}),Z(n{\scriptscriptstyle \frac{\widehat{\mu }}{2}})& Z^{}(n+{\scriptscriptstyle \frac{\widehat{\eta }}{2}}),Z(n{\scriptscriptstyle \frac{\widehat{\eta }}{2}})\mathrm{\Phi }^{(\rho \eta )}\hfill \\ \hfill \mathrm{\Psi }^{(\mu \nu )}Z^{}(n+{\scriptscriptstyle \frac{\widehat{\nu }}{2}}),Z(n{\scriptscriptstyle \frac{\widehat{\nu }}{2}})& Z^{}(n+{\scriptscriptstyle \frac{\widehat{\rho }}{2}}),Z(n{\scriptscriptstyle \frac{\widehat{\rho }}{2}})\mathrm{\Psi }^{(\rho \eta )}.\hfill \end{array}$$
(55)
In general, one of the two indices in the couple $`(\mu \nu )`$ specifies the direction of the link carrying the matrix $`Z`$ or $`Z^{}`$, while the other index shows which entries of the matrix are concerned: for instance, $`\overline{\varphi }_R^{(\mu ,\nu )}`$ couples to the entries $`Z^{}(n+{\scriptscriptstyle \frac{\widehat{\mu }}{2}})_{\nu ,.}`$, while $`\varphi _R^{(\rho \eta )}`$ couples to the entries $`Z(n{\scriptscriptstyle \frac{\widehat{\eta }}{2}})_{.,\rho }`$.
As a result, the couplings between the site quadruplets satisfy some “selection rules”, which mean that the matrix $``$ contains many $`4\times 4`$ empty blocks. The non-empty blocks connect the following pairs of quadruplets:
$$\begin{array}{cc}\hfill \mathrm{\Phi }^{(\mu \nu )}\mathrm{\Phi }^{(\rho \eta )}& \text{ iff }\mu =\eta \hfill \\ \hfill \mathrm{\Phi }^{(\mu \nu )}\mathrm{\Psi }^{(\rho \eta )}& \text{ iff }\mu =\rho \hfill \\ \hfill \mathrm{\Psi }^{(\mu \nu )}\mathrm{\Phi }^{(\rho \eta )}& \text{ iff }\nu =\eta \hfill \\ \hfill \mathrm{\Psi }^{(\mu \nu )}\mathrm{\Psi }^{(\rho \eta )}& \text{ iff }\nu =\rho .\hfill \end{array}$$
(56)
By analogy with the $`2`$-dimensional case, we will call $`V_\mu ^{\nu ,\eta }`$ the matrix coupling $`\mathrm{\Phi }^{(\mu \nu )}`$ with $`\mathrm{\Psi }^{(\mu \eta )}`$, and $`W_\nu ^{\mu ,\rho }`$ the matrix coupling $`\mathrm{\Psi }^{(\mu \nu )}`$ with $`\mathrm{\Phi }^{(\rho \nu )}`$. The structure of these matrices is similar to the ones in Eq. (29), except that the matrices $`Z`$, $`Z^{}`$ are replaced by $`2\times 2`$ submatrices:
$$\begin{array}{cc}\hfill V_\mu ^{\nu ,\eta }(n)& =\left(\begin{array}{cccc}Z_{\nu ,\eta }^{}(n+{\scriptscriptstyle \frac{\widehat{\mu }}{2}})& Z_{\nu ,\eta }^{}(n+{\scriptscriptstyle \frac{\widehat{\mu }}{2}})& 0& 0\\ Z_{\nu ,\eta }^{}(n+{\scriptscriptstyle \frac{\widehat{\mu }}{2}})& Z_{\nu ,\eta }^{}(n+{\scriptscriptstyle \frac{\widehat{\mu }}{2}})& 0& 0\\ 0& 0& Z_{\nu ,\eta }(n{\scriptscriptstyle \frac{\widehat{\mu }}{2}})& Z_{\nu ,\eta }(n{\scriptscriptstyle \frac{\widehat{\mu }}{2}})\\ 0& 0& Z_{\nu ,\eta }(n{\scriptscriptstyle \frac{\widehat{\mu }}{2}})& Z_{\nu ,\eta }(n{\scriptscriptstyle \frac{\widehat{\mu }}{2}})\end{array}\right),\hfill \\ \hfill W_\nu ^{\mu ,\rho }(n)& =\left(\begin{array}{cccc}Z_{\mu ,\rho }^{}(n+{\scriptscriptstyle \frac{\widehat{\nu }}{2}})& 0& Z_{\mu ,\rho }^{}(n+{\scriptscriptstyle \frac{\widehat{\nu }}{2}})& 0\\ 0& Z_{\mu ,\rho }(n{\scriptscriptstyle \frac{\widehat{\nu }}{2}})& 0& Z_{\mu ,\rho }(n{\scriptscriptstyle \frac{\widehat{\nu }}{2}})\\ Z_{\mu ,\rho }^{}(n+{\scriptscriptstyle \frac{\widehat{\nu }}{2}})& 0& Z_{\mu ,\rho }^{}(n+{\scriptscriptstyle \frac{\widehat{\nu }}{2}})\\ 0& Z_{\mu ,\rho }(n{\scriptscriptstyle \frac{\widehat{\nu }}{2}})& 0& Z_{\mu ,\rho }(n{\scriptscriptstyle \frac{\widehat{\nu }}{2}})\end{array}\right).\hfill \end{array}$$
(57)
We have seen that in $`d3`$, the fields $`\mathrm{\Phi }^{(\mu \nu )}`$ and $`\mathrm{\Phi }^{(\rho \mu )}`$ are also coupled, through a matrix $`X_\mu ^{\nu \rho }`$; similarly, $`\mathrm{\Psi }^{(\mu \nu )}`$ and $`\mathrm{\Psi }^{(\nu \eta )}`$ are coupled through a matrix $`Y_\nu ^{\mu \eta }`$. As for $`V`$ and $`W`$, the lower index refers to the direction of the links carrying the elements of $`Z`$, $`Z^{}`$ which appear in $`X`$ (or $`Y`$). These matrices have similar forms as the matrices $`V`$, $`W`$ above (we won’t need their explicit expression in the following). These four sets of matrices can be grouped separately into matrices of size $`N_b\times N_b`$, which we call $`𝒱(n)`$, $`𝒲(n)`$, $`𝒳(n)`$, $`𝒴(n)`$. These four matrices make up the complete coupling matrix $`(n)`$: the action reads reads
$$S_Z[n]=\left(\begin{array}{c}\mathrm{\Phi }^{(\mu \nu )}\\ \mathrm{\Psi }^{(\mu \nu )}\end{array}\right)^{}\left(\begin{array}{cc}𝒳& 𝒱\\ 𝒲& 𝒴\end{array}\right)\left(\begin{array}{c}\mathrm{\Phi }^{(\mu \nu )}\\ \mathrm{\Psi }^{(\mu \nu )}\end{array}\right).$$
(58)
Taking the mass term into account, the integral over the auxiliary fields yields
$$Det(m_b𝕀)^1\mathrm{exp}\left\{\mathrm{Tr}\mathrm{ln}(1m_b^1)\right\}\mathrm{exp}\left\{m_b^1\mathrm{Tr}+\frac{m_b^2}{2}\mathrm{Tr}^2\right\},$$
(59)
where we performed the large-$`m_b`$ expansion up to second order. To analyze the traces, we use the ‘selection rules’ given by (56). $`X_\mu ^{\nu \rho }`$ connects planes $`(\mu \nu )>(\rho \mu )`$, therefore its block appears under the diagonal in the matrix $`𝒳`$; on the opposite, $`Y_\nu ^{\mu \eta }`$ connects planes $`(\mu \nu )<(\nu \eta )`$, so its block is over the diagonal in $`𝒴`$. As a result, $`\mathrm{Tr}=0`$, and $`\mathrm{Tr}𝒳^2=\mathrm{Tr}𝒴^2=0`$. Therefore, the first nontrivial term appears at the order $`1/m_b^2`$, and takes the value $`\mathrm{Tr}^2=2\mathrm{T}\mathrm{r}(𝒱𝒲)`$.
To compute this term, we notice that the block $`V_\mu ^{\nu ,\eta }`$ connects planes $`(\mu \nu )`$, $`(\mu \eta )`$ sharing the same lower index $`\mu `$ (that is, positive generators $`e_{\mu \nu }`$, $`e_{\mu \eta }`$ situated on the same row); on the opposite, a block $`W_\nu ^{\mu ,\rho }`$ connects generators situated on the same column. Therefore, through $`𝒱`$ we can jump along a row, and through $`𝒲`$ we jump along a column. When computing $`2\mathrm{T}\mathrm{r}(𝒱𝒲)`$ we want to be back at the initial position after two jumps, so that only ‘immobile jumps’ are allowed:
$$\frac{1}{2}\mathrm{Tr}^2=\mathrm{Tr}(𝒱𝒲)=\underset{\mu <\nu }{}\mathrm{Tr}(V_\mu ^{\nu ,\nu }W_\nu ^{\mu ,\mu }).$$
(60)
Thus, to this order the planes are “decoupled” from one another, and the contribution of each plane is identical to what we had found in the two-dimensional framework (see Eq. (34)):
$$\begin{array}{c}\mathrm{Tr}\left(𝒱(n)𝒲(n)\right)=\underset{1\mu <\nu d}{}Z_{\nu ,\nu }^{}(n+{\scriptscriptstyle \frac{\widehat{\mu }}{2}})Z_{\mu ,\mu }^{}(n+{\scriptscriptstyle \frac{\widehat{\nu }}{2}})+Z_{\nu ,\nu }^{}(n+{\scriptscriptstyle \frac{\widehat{\mu }}{2}})Z_{\mu ,\mu }(n{\scriptscriptstyle \frac{\widehat{\nu }}{2}})\hfill \\ \hfill +Z_{\nu ,\nu }(n{\scriptscriptstyle \frac{\widehat{\mu }}{2}})Z_{\mu ,\mu }^{}(n+{\scriptscriptstyle \frac{\widehat{\nu }}{2}})+Z_{\nu ,\nu }(n{\scriptscriptstyle \frac{\widehat{\mu }}{2}})Z_{\mu ,\mu }(n{\scriptscriptstyle \frac{\widehat{\nu }}{2}}).\end{array}$$
(61)
Thus, as in two and three dimensions, only the diagonal elements of the $`Z`$-fields contribute to the action up to order $`1/m_b^2`$. Each of the above terms is the product of two diagonal matrix elements which self-coupled auxiliary fields carried by the same plaquette, so each term can be associated with a well-defined plaquette (we come back to this property in next section). The higher-order terms in $`1/m_b`$ are more complicated, since the matrices $`X`$, $`Y`$ and non-diagonal blocks $`V`$, $`W`$ start contributing.
To summarize, the effective action to second order in $`1/m_b`$ has the same structure in any dimension:
$$S[Z]=\underset{n}{}m_b^2\underset{\mu <\nu }{}\mathrm{Tr}(V_\mu ^{\nu ,\nu }W_\nu ^{\mu ,\mu })(n)+\underset{\alpha =1}{\overset{d}{}}\mathrm{Tr}\mathrm{ln}\left(1Z(n+{\scriptscriptstyle \frac{\widehat{\alpha }}{2}})Z^{}(n+{\scriptscriptstyle \frac{\widehat{\alpha }}{2}})\right).$$
(62)
### 3.5 Stationary point of the large-$`m_b`$ effective action
The factor $`N_c`$ in front of the action $`S[Z]`$ suggests to study the large-$`N_c`$ limit of the theory, that is look for stationary points of this action with respect to variations of the matrix fields $`Z`$, $`Z^{}`$. Since we computed the action $`S[Z]`$ up to second order in $`1/m_b`$, we will keep this approximation (62) and compute its saddle-point equations.
The variation of the quadratic terms $`\mathrm{Tr}(𝒱𝒲)`$ is easy to compute from (61): it only involves variations of diagonal elements of the matrices $`Z`$ or $`Z^{}`$. On the opposite, the variation of the second term in Eq. (62) involves all matrix elements:
$$\begin{array}{cc}\hfill \delta \mathrm{Tr}\mathrm{ln}(1ZZ^{})& =\mathrm{Tr}\left[\delta ZZ^{}(1ZZ^{})^1+\delta Z^{}Z(1Z^{}Z)^1\right]\hfill \\ & =\underset{a,b}{}\delta Z_{ab}\left(Z^{}(1ZZ^{})^1\right)_{ba}+\delta Z_{ab}^{}\left(Z(1Z^{}Z)^1\right)_{ba}.\hfill \end{array}$$
(63)
Therefore, setting $`\frac{\delta S[Z]}{\delta Z_{ab}}=0`$ for all $`ab`$ implies that the matrix $`Z(1Z^{}Z)^1`$ is diagonal; this implies that $`Z`$ is itself a diagonal matrix. We then compute the saddle point equations with respect to variations of the diagonal elements $`Z_{aa}`$. For any site $`n`$ and $`\mu \nu `$, Eq. (61) yields the following variations:
$$\frac{\delta S[Z]}{\delta Z_{\nu ,\nu }^{}(n+{\scriptscriptstyle \frac{\widehat{\mu }}{2}})}=\frac{Z_{\nu ,\nu }(n+{\scriptscriptstyle \frac{\widehat{\mu }}{2}})}{1|Z_{\nu ,\nu }(n+{\scriptscriptstyle \frac{\widehat{\mu }}{2}})|^2}m_b^2Z_{\mu ,\mu }^{}(n+{\scriptscriptstyle \frac{\widehat{\nu }}{2}}),$$
(64)
and similar expressions for the variation of $`S[Z]`$ with respect to the components $`Z_{\nu ,\nu }^{}(n+{\scriptscriptstyle \frac{\widehat{\mu }}{2}})`$, $`Z_{\nu ,\nu }(n{\scriptscriptstyle \frac{\widehat{\mu }}{2}})`$ and $`Z_{\nu ,\nu }(n{\scriptscriptstyle \frac{\widehat{\mu }}{2}})`$. Setting these variations to zero, we get the full set of saddle-points equations. These equations obviously admit the trivial configuration $`Z0`$ as solution. Taking into account the condition $`m_b1`$, one can show that this solution is the unique one.
It therefore makes sense to expand the action (62) to quadratic order in $`Z`$, $`Z^{}`$:
$$S[Z]_{\mathrm{quad}}=m_b^2\underset{n}{}\left(\underset{\mu <\nu }{}\mathrm{Tr}(V_\mu ^{\nu ,\nu }W_\nu ^{\mu ,\mu })(n)\underset{\alpha =1}{\overset{d}{}}\underset{a,b=1}{\overset{d^{}}{}}m_b^2|Z_{ab}(n+{\scriptscriptstyle \frac{\widehat{\alpha }}{2}})|^2\right).$$
(65)
From the expression (61), the above action seems to describe free bosonic fields. The nondiagonal elements $`Z_{ab}`$ with $`ab`$ are nondynamical, since they only appear in the mass term. On the opposite, the diagonal terms $`Z_{aa}`$ appear both in the mass term and the ’kinetic energy term’ (61), so they seem to correspond to propagating modes. This is actually not the case: as we already noticed, the terms Eq. (61) only couple matrix elements related to the same plaquette, so that these fields can only propagate around one plaquette. The diagonal fields are therefore non-propagating modes as well, so the above quadratic action is non-dynamical. This is not so surprising, since the auxiliary bosonic fields $`\varphi (n)`$ were from the beginning also confined to one plaquette. The same phenomenon persists if one includes several generations $`n_b>1`$.
Propagation can be induced by including higher-order terms in $`ZZ^{}`$ when expanding the logarithm. This way, one obtains a quartic contribution $`\mathrm{Tr}(ZZ^{})^2`$, which allows to couple together different diagonal elements through non-diagonal ones. This contribution includes for instance terms of the form $`Z_{2,2}Z_{2,2}^{}Z_{2,2}Z_{2,2}^{}`$ (all elements on the link $`(n+\widehat{1}/2)`$), which couple fields carried by two adjacent plaquettes, namely the plaquette $`(n,1,2)`$ carrying $`Z_{2,2}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})`$, and the plaquette $`(n,1,2)`$ carrying $`Z_{2,2}(n+{\scriptscriptstyle \frac{\widehat{1}}{2}})`$ (see Fig. (4)).
## 4 Concluding remarks
We have considered the dual formulation of the lattice theory which induces Wilson’s pure gauge action after integrating over auxiliary bosonic fields, in the limit of large mass and many “generations”. In our model the “flavor” degrees of freedom are associated not only with the number of the “generations” of the inducing field, but also with a particular plaquette; besides, we need fields with left. resp.. right “chirality”, which doubles the number of flavor degrees of freedom.
We investigate the properties of the lattice model in the simpler case of one generation. The structure of the inducing theory allows us to apply the color-flavor transformation to obtain a ‘dual’ effective theory in terms of colorless matrices $`Z`$ carried by the lattice links. After integrating over the auxiliary bosons, we obtain an effective action uniquely in terms of the $`Z`$ fields, which is computed explicitly in the limit of large auxiliary mass, leading to a trivial non-propagating theory in the large-$`N_c`$ limit.
The color-flavor transformation for the $`SU(N_c)`$ gauge group yields some differences, related with the decomposition of the colorless sector into disconnected subsectors labeled by the baryonic charge $`Q`$ . Our above derivations correspond to the sectors $`Q=0`$ with no contribution of closed baryon loops . The investigation of the effect of these loops is in progress. Note that the choice of $`U(N_c)`$ instead of the realistic $`SU(N_c)`$ is irrelevant in the large-$`N_c`$ limit.
#### Acknowledgements
This research is inspired by numerous discussions with S. Nonnenmacher, who suggested a general strategy of the integration over the auxiliary fields in $`d`$-dimensions , and with J. Budczies, who independently constructed $`d=2`$ one-plaquette effective action of the Z-field . I am grateful to D. Diakonov, V. Petrov, M. Polyakov and M.R. Zirnbauer for useful discussions and comments. I would like to acknowledge the hospitality at the Abdus Salam International Center for Theoretical Physics where this work was completed.
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# On Gauge Choice of Spherically Symmetric 3-Branes
## I Introduction
A number of current unification theories such as string theory and M-theory suggest that we may live in a world that has more than three spatial dimensions. Because only three of these are presently observable, one has to explain why the others are hidden from detection. One such explanation is the so-called Kaluza-Klein (KK) compactification, according to which the size of the extra dimensions is very small (often taken to be on the order of the Planck length). As a consequence, modes that have momentum in the directions of the extra dimensions are excited at currently inaccessible energies.
Recently, the braneworld scenario has dramatically changed this point of view and, in the process, received a great deal of attention. At present, there are a number of proposed models (See, for example, and references therein.). In particular, Arkani-Hamed et al (ADD) pointed out that the extra dimensions need not necessarily be small and may even be on the scale of millimeters. This model assumes that Standard Model fields are confined to a three (spatial) dimensional hypersurface (a 3-brane) living in a larger dimensional bulk space while the gravitational field propagates in the whole bulk. Additional fields may live only on the brane or in the whole bulk, provided that their current undetectability is consistent with experimental bounds. One of the most attractive features of this model is that it may potentially resolve the long standing hierarchy problem, namely the large difference in magnitudes between the Planck and electroweak scales.
In a different model, Randall and Sundrum (RS) showed that if the self-gravity of the brane is included, gravitational effects can be localized near the brane at low energy and the 4D newtonian gravity will be reproduced on the brane even in the presence of infinitely large extra dimensions. In this model, the $`4D`$ Planck scale, $`M_{Pl}`$, is determined by the curvature of the extra dimensions rather than by their size, as proposed in .
The RS model <sup>*</sup><sup>*</sup>*In this Letter we are mainly concerned with the so-called RS2 model, in which only one brane exists. was soon generalized to include arbitrary matter fields on the brane . In particular, Shiromizu, Maeda and Sasaki (SMS) considered the embedding of a 3-brane ($`M,{}_{}{}^{(4)}g_{AB}^{}`$) in a 5D bulk ($`V,{}_{}{}^{(5)}g_{AB}^{}`$), where the 3-brane metric is given by $`{}_{}{}^{(4)}g_{AB}^{}={}_{}{}^{(5)}g_{AB}^{}n_An_B`$, and $`n_A`$ is the unit normal vector to the brane. Note that we use Greek indices to run from $`0`$ to $`3`$, uppercase Latin indices to run from $`0`$ to $`D1`$, and lowercase Latin indices to run from $`0`$ to $`D2`$. Using the Gauss–Codacci relations, SMS wrote the 5D Einstein field equations $`{}_{}{}^{(5)}G_{AB}^{}=\kappa _5^2{}_{}{}^{(5)}T_{AB}^{}`$ in the form
$`{}_{}{}^{(4)}R_{AB}^{}{\displaystyle \frac{1}{2}}{}_{}{}^{(4)}g_{AB}^{}{}_{}{}^{(4)}R={}_{}{}^{(4)}𝒯_{AB}^{},`$ (1.1)
$`D_CK_A^CD_AK=\kappa _5^2{}_{}{}^{(5)}T_{BC}^{}{}_{}{}^{(4)}g_{A}^{C}n^B,`$ (1.2)
where these equations are understood to apply in each of the two regions, $`V^+(z0)`$ and $`V^{}(z0)`$, and $`z`$ denotes the coordinate of the extra dimension and $`z=0`$ is the location of the brane. The quantity $`{}_{}{}^{(4)}𝒯_{AB}^{}`$ is given by
$`{}_{}{}^{(4)}𝒯_{AB}^{}`$ $``$ $`{\displaystyle \frac{2\kappa _{5}^{}{}_{}{}^{2}}{3}}\left\{{}_{}{}^{(5)}T_{CD}^{}{}_{}{}^{(4)}g_{A}^{C}{}_{}{}^{(4)}g_{B}^{D}+\left[{}_{}{}^{(5)}T_{CD}^{}n^Cn^D{\displaystyle \frac{1}{4}}{}_{}{}^{(5)}T_{C}^{C}\right]{}_{}{}^{(4)}g_{AB}^{}\right\}`$ (1.4)
$`+KK_{AB}K_A^CK_{BC}{\displaystyle \frac{1}{2}}{}_{}{}^{(4)}g_{AB}^{}\left(K^2K^{CD}K_{CD}\right)_{AB},`$
$`_{AB}`$ $``$ $`{}_{}{}^{(5)}C_{FCD}^{E}n_En^C{}_{}{}^{(4)}g_{A}^{F}{}_{}{}^{(4)}g_{B}^{D},`$ (1.5)
and $`{}_{}{}^{(5)}C_{BCD}^{A}`$ denotes the Weyl tensor of the bulk.
The boundary conditions at $`z=0`$ are simply the Israel junction conditions ,
$$\left[K_{AB}\right]^{}=\kappa _5^2\left(S_{AB}\frac{1}{3}{}_{}{}^{(4)}g_{AB}^{}S\right),$$
(1.6)
where
$`\left[K_{AB}\right]^{}`$ $``$ $`\underset{z0^+}{lim}K_{AB}^+\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\underset{z0^{}}{lim}K_{AB}^{},`$ (1.7)
$`S_{AB}`$ $``$ $`{}_{}{}^{(4)}T_{AB}^{}\lambda {}_{}{}^{(4)}g_{AB}^{},`$ (1.8)
with $`\lambda `$ and $`{}_{}{}^{(4)}T_{AB}^{}`$ being, respectively, the cosmological constant and the energy-momentum tensor on the 3-brane. Combining the assumption of $`Z_2`$ symmetry with Eqs. (1.6) and (1.1) in the limit $`z0^\pm `$, SMS obtained the effective 4D Einstein field equations on the 3-brane
$${}_{}{}^{(4)}G_{AB}^{}=\mathrm{\Lambda }_4{}_{}{}^{(4)}g_{AB}^{}\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }+8\pi G_4{}_{}{}^{(4)}T_{AB}^{}\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }+\kappa _5^4\pi _{AB}\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }_{AB},$$
(1.9)
where $`G_4\kappa _{5}^{}{}_{}{}^{4}\lambda /(48\pi )`$. These equations include two correction terms to the original Einstein equations, namely $`_{AB}`$ and $`\pi _{AB}`$. In the weak field limit, the first term gives rise to the massive KK modes of the graviton . The second term, $`\pi _{AB}`$, is negligible for weak fields. However, in strong field situations, such as those at the threshold of black hole formation, it is expected to play a significant role. Indeed, it would appear to be this term which is the origin of the result of that, given the 4D projected Einstein equations (1.9) and the Israel matching conditions (1.6), the exterior of a collapsing homogeneous dust cloud cannot be static. This result was further generalized to other cases , and represents a significant departure from the familiar result in Einstein’s theory, in which the vacuum exterior must be the static Schwarschild solution. A number of other recent results suggest additional phenomena different from the standard predictions of general relativity . In particular, it was argued that static braneworld black holes might not exist at all .
Braneworld scenarios have further been promoted by the possibility that they may provide the origin of dark energy , which is needed in order to explain why our universe is currently accelerating .
In this Letter, we report on some results concerning the gauge choice for matter fields confined on a spherically symmetric 3-brane. We show that the boundary conditions (1.6) serve as very strong restrictions on the possible dependence of the matter fields on the spacetime coordinates. In particular, for a particular choice of the gauge, a scalar field or a Yang-Mills field can be only either time-dependent or radial-coordinate dependent, while for a perfect fluid its radial velocity must vanish. Moreover, these conclusions would appear to be true not only for the generalized RS models in a 5D bulk, but also for 3-branes in higher dimensional spacetimes . This is quite different from its four-dimensional counterpart.
Before showing these results, let us first give a brief review on the gauge choice of a 4D spacetime with spherical symmetry, for which the general metric can be cast, without loss of generality, in the form,
$$ds_4^2=g_{ab}\left(x^c\right)dx^adx^b+s^2\left(x^c\right)d\mathrm{\Omega }^2,(a,b,c=0,1),$$
(1.10)
where $`d\mathrm{\Omega }^2d\theta ^2+\mathrm{sin}^2\theta d\phi ^2`$. Clearly, the form of the metric is invariant under the coordinate transformations,
$$x^0=x^0(x_{}^{}{}_{}{}^{0},x_{}^{}{}_{}{}^{1}),x^1=x^1(x_{}^{}{}_{}{}^{0},x_{}^{}{}_{}{}^{1}).$$
(1.11)
Using these two degrees of the gauge freedom, one can choose different gauges for different matter fields. Because of the complexity of the Einstein field equations, such choices often are crucial in studying the problem. For example, for a collapsing perfect fluid, one usually chooses the so-called comoving gauge,
$$g_{01}\left(x^c\right)=0,u_A=\left(g_{00}\right)^{1/2}\delta _A^0,$$
(1.12)
so that the metric takes the form,
$$ds_4^2=g_{00}\left(x^c\right)\left(dx^0\right)^2+g_{11}\left(x^c\right)\left(dx^1\right)^2+s^2\left(x^c\right)d\mathrm{\Omega }^2,(c=0,1),$$
(1.13)
where $`u_A`$ denotes the four-velocity of the fluid, and $`x^0`$ is the time-like coordinate. For a collapsing scalar field, on the othe hand, a possible choice of the gauge is
$$g_{01}\left(x^c\right)=0,s\left(x^c\right)=x^1r,$$
(1.14)
so that the metric takes the form,
$$ds_4^2=g_{00}(t,r)dt^2+g_{11}(t,r)dr^2+r^2d\mathrm{\Omega }^2,$$
(1.15)
where $`tx^0`$ and $`\varphi =\varphi (t,r)`$, with $`\varphi `$ denoting the scalar field. Certainly, depending on a specific problem to be considered, other gauges can be chosen.
## II Gauge Freedom and Gauge Choice For a Spherical 3-Brane
To begin with, consider the general action describing a 3-brane embedded in a D–dimensional bulk ,
$`S`$ $`=`$ $`{\displaystyle \frac{1}{2\kappa _D^2}}{\displaystyle _{M_D}}d^Dx\sqrt{{}_{}{}^{(D)}g}\left({}_{}{}^{(D)}R2\mathrm{\Lambda }_D+2\kappa _D^2_m^B\right)`$ (2.2)
$`+{\displaystyle \frac{1}{2\kappa ^2}}{\displaystyle _{M_4}}d^4x\sqrt{g}\left(2\mathrm{\Lambda }+2\kappa ^2_m\right),`$
where $`{}_{}{}^{(D)}R,\mathrm{\Lambda }_D`$ and $`_m^B`$ ($`R,\mathrm{\Lambda },_m`$) denote, respectively, the Ricci scalar, cosmological constant and matter content of the bulk (of the 3-brane). The constants $`\kappa _D`$ and $`\kappa `$ are related to the Planck scales $`M`$ and $`M_{Pl}`$, respectively, by
$`\kappa _D^2`$ $`=`$ $`8\pi {}_{}{}^{(D)}G=M^{2D},`$ (2.3)
$`\kappa ^2`$ $`=`$ $`8\pi G=M_{Pl}^2,`$ (2.4)
with $`Dn+4`$. The bulk metric is $`{}_{}{}^{(D)}g_{MN}^{}`$ and $`g_{\mu \nu }`$ denotes the induced metric on the 3-brane, located on the surface $`\mathrm{\Phi }(x^A)=0`$. Taking a different approach from SMS, we vary Eq. (2.2) with respect to $`{}_{}{}^{(D)}g_{MN}^{}`$ and $`g_{\mu \nu }`$ to get the full D–dimensional Einstein field equations in the form
$`{}_{}{}^{(D)}G_{MN}^{+}`$ $`=`$ $`\kappa _D^2T_{MN}^{B+}\mathrm{ }\mathrm{\Lambda }_D{}_{}{}^{(D)}g_{MN}^{+},(\mathrm{\Phi }0),`$ (2.5)
$`{}_{}{}^{(D)}G_{MN}^{}`$ $`=`$ $`\kappa _D^2T_{MN}^B\mathrm{ }\mathrm{\Lambda }_D{}_{}{}^{(D)}g_{MN}^{},(\mathrm{\Phi }0),`$ (2.6)
$`{}_{}{}^{(4)}G_{\mu \nu }^{Im}`$ $`=`$ $`\kappa _D^2\left(T_{\mu \nu }{\displaystyle \frac{\mathrm{\Lambda }}{\kappa ^2}}g_{\mu \nu }\right),(\mathrm{\Phi }=0),`$ (2.7)
where $`T_{MN}^B`$ and $`T_{\mu \nu }`$ denote the stress energy tensor of the bulk and of the 3-brane, respectively. The quantities with superscript “$`+`$” (“$``$”) denote those calculated in the region $`\mathrm{\Phi }0`$ ($`\mathrm{\Phi }0`$), and $`{}_{}{}^{(4)}G_{\mu \nu }^{Im}`$ denotes the delta-function-like (impulsive) part of $`G_{MN}`$ with support on the 3-brane,
$${}_{}{}^{(D)}G_{MN}^{}={}_{}{}^{(D)}G_{MN}^{+}H\left(\mathrm{\Phi }\right)+{}_{}{}^{(D)}G_{MN}^{}\left[1H\left(\mathrm{\Phi }\right)\right]+{}_{}{}^{(4)}G_{\mu \nu }^{Im}\delta _M^\mu \delta _N^\nu \delta \left(\mathrm{\Phi }\right),$$
(2.8)
where $`\delta \left(\mathrm{\Phi }\right)`$ denotes the Dirac delta function, and $`H(\mathrm{\Phi })`$ the Heavside function, defined as
$$H\left(\mathrm{\Phi }\right)=\{\begin{array}{cc}1,\hfill & \mathrm{\Phi }0\text{,}\hfill \\ 0,\hfill & \mathrm{\Phi }<0\text{.}\hfill \end{array}$$
(2.9)
Eq.(2.8) can easily be obtained by the following considerations. Let us first denote the region with $`\mathrm{\Phi }0`$ as $`V^+`$, the region with $`\mathrm{\Phi }0`$ as $`V^{}`$, and the hypersurface $`\mathrm{\Phi }=0`$ as $`\mathrm{\Sigma }`$. Then, since the Einstein field equations involve the second-order derivatives of the metric coefficients, one can see that the metric must be at least $`C^2`$ in regions $`V^\pm `$ and $`C^0`$ across the hypersurface $`\mathrm{\Phi }=0`$, so that the Einstein field equations (or any second-order differential equations involved) hold in the sense of distributions . Consequently, the metric $`g_{AB}`$ in the whole spacetime can be written as
$$g_{AB}=g_{AB}^+H\left(\mathrm{\Phi }\right)+g_{AB}^{}\left[1H\left(\mathrm{\Phi }\right)\right],$$
(2.10)
where quantities with superscripts $`\mathrm{`}\mathrm{`}+\mathrm{"}`$ ($`\mathrm{`}\mathrm{`}\mathrm{"}`$) denote the ones defined in $`V^+`$ ($`V^{}`$). Hence, we find that
$`g_{AB,C}`$ $`=`$ $`g_{AB,C}^+H\left(\mathrm{\Phi }\right)+g_{AB,C}^{}\left[1H\left(\mathrm{\Phi }\right)\right],`$ (2.11)
$`g_{AB,CD}`$ $`=`$ $`g_{AB,CD}^+H\left(\mathrm{\Phi }\right)+g_{AB,CD}^{}\left[1H\left(\mathrm{\Phi }\right)\right]+\left[g_{AB,C}\right]^{}\mathrm{\Phi }_{,D}\delta \left(\mathrm{\Phi }\right),`$ (2.12)
where $`()_{,C}()/x^C`$ and
$$\left[g_{AB,C}\right]^{}\underset{\mathrm{\Phi }0^+}{lim}\frac{g_{AB}^+}{x^C}\underset{\mathrm{\Phi }0^{}}{lim}\frac{g_{AB}^{}}{x^C}.$$
(2.14)
From Eq.(2.11) and the following,
$`H^m\left(\mathrm{\Phi }\right)=H\left(\mathrm{\Phi }\right),\left[1H\left(\mathrm{\Phi }\right)\right]^m=1H\left(\mathrm{\Phi }\right),`$ (2.15)
$`H\left(\mathrm{\Phi }\right)\left[1H\left(\mathrm{\Phi }\right)\right]=0,\left[1H\left(\mathrm{\Phi }\right)\right]\delta (\mathrm{\Phi })={\displaystyle \frac{1}{2}}\delta (\mathrm{\Phi })=H\left(\mathrm{\Phi }\right)\delta (\mathrm{\Phi }),`$ (2.16)
where $`m`$ is an integer, we can easily see that the Einstein tensor $`G_{AB}`$ can be written, in general, in the form of Eq.(2.8).
In the 5D case a 3-brane is a hypersurface of a 5D bulk, and one can show that Eq. (2.7) is identical to Eq. (1.6) when written out in terms of the extrinsic curvature of the 3-brane . Similarly, using the Gauss–Codacci relations one can show that Eqs. (1.1) and (1.2) follow from Eqs. (2.5) and (2.6). It is, at the same time, worth emphasizing that Eqs. (2.52.6) contain additional information not present in Eqs. (1.11.2). This includes, for instance, information about the evolution of $`_𝒜`$. For the case that $`D6`$, a 3-brane is a surface of co–dimension $`(D4)`$ with respect to the bulk, and the problem becomes more subtle. In particular, the generalization of the Gauss–Codacci relations and the Israel junction conditions to these cases has not, to our knowledge, been worked out. For this reason, in the $`D6`$ cases we will only consider models with additional symmetries. In this way, a 3-brane can be considered as a degenerate hypersurface. Indeed, this assumption turns out to include most of the braneworld models with higher dimensional bulks which have been studied so far .
### A $`D=5`$
Considering first the case $`D=5`$, the most general bulk metric with a $`S^2`$ symmetry takes the form,
$$ds^2=g_{ij}\left(x^k\right)dx^idx^j\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }+s^2\left(x^k\right)d\mathrm{\Omega }^2,$$
(2.17)
where $`i`$, $`j`$ and $`k`$ are taken here to range over $`0`$, $`\mathrm{\hspace{0.17em}1}`$ and $`2`$ with $`zx^2`$. The location of a 3-brane with spherical symmetry in general can be written as
$$\mathrm{\Phi }(x^0,x^1,z)=0.$$
(2.18)
Note that the form of the metric (2.17) is invariant under the coordinate transformations,
$$x^i=\mathrm{ }f^i\left(\overline{x}^j\right),(i,j=0,1,2).$$
(2.19)
As a result, using these three degrees of freedom, we can choose coordinates such that the brane is always located on the hypersurface $`z=0`$ and
$$\mathrm{\Phi }(x^0,x^1,z)=z,g_{0z}(x^0,x^1,z)=g_{1z}(x^0,x^1,z)=0.$$
(2.20)
This choice of coordinates will be referred to as the canonical gauge.
Using this form for the metric together with the definition of the extrinsic curvature, $`K_{AB}h_A{}_{}{}^{C}h_{B}^{}{}_{}{}^{D}_{C}^{}n_D`$, we find that
$$K_{\mu \nu }=(2N)^1\frac{g_{\mu \nu }(x^\alpha ,z)}{z},$$
(2.21)
where $`n_A=N\delta _A^z,N\sqrt{g_{zz}}`$, and $`_C`$ denotes the covariant derivative with respect to $`{}_{}{}^{(D)}g_{AB}^{}`$. For the metric (2.17) in the canonical gauge we find that the Israel junction conditions, (1.6), yield
$$T_{\mu \nu }=\mathrm{ }\mathrm{ }\lambda g_{\mu \nu }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\frac{1}{2\kappa _5^2N}\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\left(\mathrm{ }\mathrm{ }\left[g_{\mu \nu ,z}\right]^{}\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }g_{\mu \nu }g^{\alpha \beta }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\left[g_{\alpha \beta ,z}\right]^{}\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\right).$$
(2.22)
It should be noted that even in the canonical gauge, there is residual coordinate freedom on the 3-brane:
$$x^0=F^0(\overline{x}^0,\overline{x}^1),x^1=F^1(\overline{x}^0,\overline{x}^1).$$
(2.23)
We can exploit this remaining freedom and set
$$g_{01}(x^0,x^1,0)=g_{01,z}(x^0,x^1,0)=0,$$
(2.24)
so that the reduced metric on the 3-brane takes the form,
$$ds^2|_{z=0}=\gamma _{00}(x^0,x^1)\left(dx^0\right)^2+\gamma _{11}(x^0,x^1)\left(dx^1\right)^2+\gamma _{22}(x^0,x^1)d\mathrm{\Omega }^2,$$
(2.25)
where $`\gamma _{ab}(x^0,x^1)g_{ab}(x^0,x^1,0)`$. It is remarkable to note that Eq.(2.24) is possible only in the cases where one of the three energy conditions, weak, strong and dominant, holds . To show this, following Chandrasekhar , we first transform to the coordinates $`\overline{t}`$ and $`\overline{r}`$ using the transformations $`x^0=\varphi (\overline{t},\overline{r})`$ and $`x^1=\psi (\overline{t},\overline{r})`$, in which the metric takes the form
$$ds^2=bd\overline{t}^2+2cd\overline{t}d\overline{r}+dd\overline{r}^2\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }+s^2d\mathrm{\Omega }^2+N^2dz^2,$$
(2.26)
where $`b`$, $`c`$, $`d`$, and $`s`$ are functions of $`\overline{t},\overline{r}`$ and $`z`$, and have the properties ,
$$b(\overline{t},\overline{r},0)=d(\overline{t},\overline{r},0),c(\overline{t},\overline{r},0)=0,$$
(2.27)
at $`z=0`$. If we now make another coordinate transformation $`\overline{t}=\mathrm{\Phi }(t,r),\overline{r}=\mathrm{\Psi }(t,r)`$, it is straightforward to show that the conditions $`g_{tr}(t,r,0)=0`$ and $`g_{tr,z}(t,r,0)=0`$ reduce to
$`\mathrm{\Phi }_{,t}=A\mathrm{\Psi }_{,t},`$ (2.28)
$`\mathrm{\Phi }_{,r}\mathrm{ }\mathrm{ }=\mathrm{ }\mathrm{ }A^1\mathrm{\Psi }_{,r},`$ (2.29)
$`(A\mathrm{\Psi }_{,t})_{,r}=(A^1\mathrm{\Psi }_{,r})_{,t},`$ (2.30)
where
$`A(t,r)`$ $``$ $`{\displaystyle \frac{1}{2c_{,z}}}\left[(b_{,z}d_{,z})\pm \mathrm{\Delta }^{1/2}\right]_{z=0},`$ (2.31)
$`\mathrm{\Delta }(t,r)`$ $``$ $`(b_{,z}d_{,z})^24c_{,z}{}_{}{}^{2}|_{z=0}^{}.`$ (2.32)
Eq.(2.30) represents the integrability condition of Eqs.(2.28) and (2.29). Since the metric coefficients is at least $`C^2`$ in regions $`V^\pm `$ and $`C^0`$ across the hypersurface $`z=0`$, we can see that, with respect to $`t`$ and $`r`$, the metric is also at least $`C^2`$ even across the hypersurface $`z=0`$. For such a $`C^2`$ metric, the theorems given in show that there will always exist a region of the $`(t,r)`$-plane for which Eqs.(2.28)-(2.30) have solutions. However, such solutions will be real only if
$$\mathrm{\Delta }0.$$
(2.33)
This condition is ensured by assuming any of the standard energy conditions . Indeed, using the Israel junction conditions (1.6), we find
$$\mathrm{\Delta }=(\kappa _5bN)^2\left[(\rho +p_{\overline{r}})^24q^2\right],$$
(2.34)
where
$$\rho \frac{T_{\overline{t}\overline{t}}}{b},p_{\overline{r}}\frac{T_{\overline{r}\overline{r}}}{d},q\frac{T_{\overline{t}\overline{r}}}{(bd)^{1/2}}.\mathrm{ }$$
(2.35)
As shown in , a necessary condition for any of the weak, dominant, and strong energy conditions to hold is $`\mathrm{\Delta }0`$. We thus conclude that a coordinate transformation exists such that (2.24) holds, provided that one of the three energy conditions holds.
As a consequence of our coordinates and Eq. (2.22), $`T_{\mu \nu }`$ must be diagonal. In particular, we have
$$T_{tr}(t,r)=0.$$
(2.36)
This represents a very strong restriction on the dependence of matter fields confined to the 3-brane on the spacetime coordinates. To see this clearly, let us first consider a scalar field $`\varphi `$, for which the stress tensor is given by
$$T_{\mu \nu }^\varphi =_\mu \varphi _\nu \varphi \frac{1}{2}g_{\mu \nu }[(\varphi )^2+2V(\varphi )].$$
(2.37)
In this case, we have
$$\mathrm{\Delta }=b^2(\varphi _t{}_{}{}^{2}\varphi _r{}_{}{}^{2})^2,T_{tr}^\varphi (t,r)=\varphi _{,t}(t,r)\varphi _{,r}(t,r).$$
(2.38)
Then, Eq.(2.36) implies
$$(i)\varphi (t,r)=\varphi (t),\text{or}(ii)\varphi (t,r)=\varphi (r).$$
(2.39)
One can show that this is also true for a spherically symmetric $`SU(2)`$ Yang-Mills field. In that case the relevant term is
$$T_{tr}^{\mathrm{YM}}w_{,t}(t,r)w_{,r}(t,r),$$
(2.40)
where $`w`$ is the Yang-Mills potential .
Similarly, if one considers a perfect fluid with stress tensor
$$T_{\mu \nu }^{\mathrm{fl}}=(\rho +p)u_\mu u_\nu +pg_{\mu \nu },$$
(2.41)
for which
$$\mathrm{\Delta }=(\rho +p)^2,T_{tr}^{\mathrm{fl}}=(\rho +p)u_tu_r,$$
(2.42)
it turns out that the fluid cannot have a radial velocity, that is, for the present choice of the gauge, we must have
$$u_r=0.$$
(2.43)
It is interesting to note that the condition (2.33) is satisfied for the cosmological constant, for which the corresponding energy-momentum tensor is given by Eq.(2.41) with $`p=\rho =\lambda `$.
The above results hold not only for $`D=5`$ but also for $`D6`$. To see this, in the following let us first consider the case where the 3-brane is a surface of co–dimension two, that is, $`D=6`$. Then, we shall further generalize our results to the case $`D>6`$.
### B $`D=6`$
Because we now have two extra spatial directions, the generalization of the RS model to include matter fields on the 3-brane becomes non-trivial. Following Israel , we will assume a cylindrical symmetry in these extra dimensions and that the 3-brane is located on the 4-dimensional surface $`\rho =0`$ where $`\rho `$ and $`\psi `$ are chosen as polar-like coordinates for the two extra dimensions, and $`\rho =0`$ is the symmetry axis for the cylindrical symmetry. This includes most of the braneworld models in six dimensional spacetimes .
Under these assumptions, it can be shown that the general bulk metric with a spherical 3-brane can be cast in the form
$`ds^2`$ $`=`$ $`\alpha ^2dt^2+a^2(dr+\beta dt)^2+s^2d\mathrm{\Omega }^2`$ (2.45)
$`+N^2d\rho ^2+f^2\rho ^2\left(d\psi +\omega dt\right)^2+g^2\rho ^2d\psi ^2,`$
where $`\omega `$ represents the rotation of the 3-brane, and all the metric coefficients are functions of $`t,r`$ and $`\rho `$, subject to the gauge (2.24). Because the symmetry axis now represents a 3-brane, certain conditions must be imposed there . For the present purpose, it is sufficient to assume that: (a) the symmetry axis must exist, that is,
$$X_\psi O(\rho ^2),$$
(2.46)
as $`\rho 0`$, where $`_\psi `$ is the cylindrical Killing vector with closed orbits, and that (b) the spacetime is free of curvature singularities on the axis, which can be assured by assuming the local flatness condition,
$$\underset{\rho 0^+}{lim}\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\mathrm{ }\frac{X_{,A}X_{,B}{}_{}{}^{(D)}g_{}^{AB}}{4X}=1.$$
(2.47)
To generalize Israel’s method to this case, we first calculate the extrinsic curvature $`K_{ab}`$ of the hypersurface $`\rho =ϵ`$ and then take the limit $`ϵ0`$. Introducing the quantities $`𝒦_{ab}`$ by
$$𝒦_{ab}\underset{\rho 0^+}{lim}\left(\sqrt{{}_{}{}^{(5)}g}K_{ab}\right),$$
(2.48)
the surface stress energy tensor can be defined as
$$T_b^a=\left(𝒦_b^a\delta _b^a𝒦_c^c\right),$$
(2.49)
provided $`𝒦_{[c}^a𝒦_{d]}^b0`$ and where we have let $`a,b=0,\mathrm{},4`$ and $`x^5=\rho `$. In this case, it can be shown that the extrinsic curvature of the hypersurface $`\rho =ϵ`$ is given by
$$K_{ab}=\frac{1}{2N}\frac{g_{ab}(x^c,ϵ)}{\rho }.$$
(2.50)
For the case $`𝒦_{[c}^a𝒦_{d]}^b=0`$, the surface stress energy tensor is instead defined as ,
$$T_b^a=\frac{2\pi }{N^2\sqrt{f^2+g^2}}\left(\sqrt{f^2+g^2}N\right)\delta _\mu ^a\delta _b^\nu ,$$
(2.51)
where $`\mu ,\nu =0,\mathrm{}3`$. In passing, we note this case also corresponds to a cosmic string in a 6D bulk .
However, in each of the above two cases it can be seen that the condition Eq.(2.36) holds. Therefore, the junction conditions across the 3-brane in a 6D bulk yield the same restrictions on the dependence of the spacetime coordinates of the matter fields confined on the 3-brane as those in the 5D case.
### C $`D>6`$
When $`D>6`$, we consider only the case where the extra $`n`$–dimensional space has an $`SO(n1)`$ symmetry so that the bulk metric can be written in the form,
$`ds^2`$ $`=`$ $`\alpha ^2dt^2+a^2(dr+\beta dt)^2+s^2d\mathrm{\Omega }^2`$ (2.53)
$`+N^2\left(d\rho ^2+\rho ^2d\mathrm{\Omega }_{n1}^2\right),`$
where $`d\mathrm{\Omega }_{n1}^2`$ denotes the metric of a unit $`(n1)`$–dimensional sphere, and all the metric coefficients are functions of $`t,r`$ and $`\rho `$. Using, as before, the coordinate freedom $`t=F_1(t^{},r^{})`$ and $`r=F_2(t^{},r^{})`$ we can always assume that Eq. (2.24) holds. From Eq. (2.53) we note that $`\rho =0`$ represents a four-dimensional spacetime with spherical symmetry. This we shall take as our 3-brane. Certainly, this is acceptable only after some (physical and or geometrical) conditions are satisfied at $`\rho =0`$. These will include, as before, that the spacetime is free of curvature singularities there. In order to generalize Israel’s method to this case, we also require that the limit
$$𝒦_{ab}=\underset{\rho 0^+}{lim}\left(\sqrt{{}_{}{}^{(D1)}g}K_{ab}\right),$$
(2.54)
exists, where $`K_{ab}`$ denotes the extrinsic curvature of the hypersurface $`\rho =ϵ`$, but now with $`a,b=0,1,\mathrm{},D2`$, and $`x^{D1}=\rho `$. With these conditions, we can define the surface stress energy tensor as that given by Eqs. (2.49) and (2.51). For the metric (2.53), it can be shown that the extrinsic curvature $`K_{ab}`$ of the hypersurface $`\rho =ϵ`$ is also given by Eq. (2.50). Substituting it into Eq. (2.49), we find again that the component $`T_{tr}(t,r)`$ vanishes identically for both of the cases described by Eqs. (2.49) and (2.51). Thus, the same restrictions on the dependence of the spacetime coordinates of the matter fields confined to the 3-brane that occur in the 5D case continue to hold for a D–dimensional bulk given by metric (2.53).
## III Conclusions
In summary, we have studied the embedding of a spherically symmetric 3-brane into a D–dimensional bulk with arbitrary matter fields both on the brane and in the bulk in the context of the braneworld scenario. We have found that, for a particular choice of gauge, the boundary (Israel’s junction) conditions across the brane together with imposition of the weak energy condition provide very strong restrictions on the dependence of matter fields confined to the 3-brane on the spacetime coordinates. As examples, a scalar field or a Yang-Mills field can be only either time-dependent or radial-coordinate dependent, while for a perfect fluid its radial velocity must vanish.
In this paper, we have studied only the function dependence of the metric coefficients for the canonical gauge. It would be very interesting to study the effects of the matter fields in the bulk on the effective 4-dimensional energy-momentum tensor $`{}_{}{}^{(4)}T_{AB}^{}`$, a subject that is under our current investigation. Another interesting problem is the applications of the results obtained in this paper to cosmology .
## Acknowledgments
The author would like to thank Rong-Gen Cai, Andrew Chamblin, Michael Christensen, Roy Maartens, and Zhong-Chao Wu for useful conversations and comments. His special thanks goes to E.W. Hirschmann for his valuable collaboration in the early stage of this work. Part of the work was done when the author was visiting Physics Department of Brigham Young University (BYU), and the Astrophysics Center, Zhejiang University of Technology (ZUT). He would like to express his gratitude to them for their hospitality. The financial assistance from Baylor University for the 2005 summer sabbatical leave, BYU, and ZUT is gratefully acknowledged.
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# 1 Introduction
## 1 Introduction
The random matrix (RM) theory was introduced originally as an approximation theory of statistics of nuclear energy levels . It should be noted that at the same time as the standard theory was established for three ensembles called the Gaussian unitary, orthogonal, and symplectic ensembles (GUE, GOE, GSE) , Dyson proposed to study such stochastic processes of interacting particles that the eigenvalue statistics of RMs are realized in distribution of particle positions on $``$ . Dyson’s Brownian motion model is a one-parameter family of $`N`$-particle systems, $`𝐙^{(\beta )}(t)=(Z_1^{(\beta )}(t),Z_2^{(\beta )}(t),\mathrm{},Z_N^{(\beta )}(t))`$, described by the stochastic differential equations
$$dZ_i^{(\beta )}(t)=dB_i(t)+\frac{\beta }{2}\underset{1jN,ji}{}\frac{1}{Z_i^{(\beta )}(t)Z_j^{(\beta )}(t)}dt,t[0,\mathrm{}),1iN,$$
(1.1)
where $`B_i(t),i=1,2,\mathrm{},N`$ are independent standard Brownian motions and the parameter $`\beta `$ equals 2, 1 and 4 for GUE, GOE and GSE, respectively. Due to the strong repulsive forces, which are long-ranged and act between any pair of particles, intersections of particle trajectories are prohibited for $`\beta 1`$ (see also ). In this one-parameter family, the $`\beta =2`$ case (i.e. the GUE case) is the simplest and the most-understood, since its equivalence with the $`N`$ particle systems of Brownian motions conditioned never to collide with each other can be proved .
The standard (Wigner-Dyson) theory has been extended by adding three chiral versions of RM ensembles in the particle physics of QCD , and by introducing the four additional ensembles so-called the Bogoliubov-de Gennes classes in the mesoscopic physics . Here we note that the chiral ensembles have a parameter $`\nu \{0,1,2,\mathrm{}\}`$ in addition to $`\beta `$. In these totally ten ensembles , chiral GUE (chGUE), class C and class D can be regarded as natural extensions of the GUE, in the sense that these eigenvalue statistics are also realized in appropriate non-colliding systems of stochastic particle systems: König and O’Connell showed that the chGUE with the parameter $`\nu \{0,1,2,\mathrm{}\}`$ corresponds to the non-colliding systems of $`2(\nu +1)`$-dimensional squared Bessel processes . The present authors clarified that the eigenvalue statistics in the classes C and D are realized by the non-colliding systems of the Brownian motions with an absorbing wall at the origin and of the Brownian motions reflecting at the origin . Since the absorbing and reflecting Brownian motions are directly related with the three-dimensional and one-dimensional Bessel processes, respectively (see, for example, ), the stochastic differential equations of these non-colliding particle systems are generally given by
$`d\stackrel{~}{Z}_i^{(\nu )}(t)=dB_i(t)+\left[{\displaystyle \frac{2\nu +1}{2}}{\displaystyle \frac{1}{\stackrel{~}{Z}_i^{(\nu )}(t)}}+{\displaystyle \underset{1jN,ji}{}}\left\{{\displaystyle \frac{1}{\stackrel{~}{Z}_i^{(\nu )}(t)\stackrel{~}{Z}_j^{(\nu )}(t)}}+{\displaystyle \frac{1}{\stackrel{~}{Z}_i^{(\nu )}(t)+\stackrel{~}{Z}_j^{(\nu )}(t)}}\right\}\right]dt,`$
$`t[0,\mathrm{}),1iN,`$ (1.2)
with reflecting barrier condition at the origin in case $`\nu =1/2`$. Therefore, the difference of (non-standard) RM ensembles can be attributed to the difference of dimensionality of the Bessel processes, whose non-colliding sets realize the statistics of the RM ensembles . Here we remind that the $`d`$-dimensional Bessel process is defined as the process of the radial coordinate (the modulus) of a Brownian motion in $`^d`$. To realize other $`104=6`$ RM ensembles by conditioned stochastic processes may be much more difficult (see ), but we demonstrated that, if we consider appropriate non-colliding systems of temporally inhomogeneous processes defined only in a finite time-interval $`[0,T]`$, we can observe the transitions of distributions into the 6 distributions as the time $`t`$ approaches the final time $`T`$ . The interesting fact is that the processes that can be used instead of the Bessel processes (1.2) should have one more parameter $`\kappa `$ in addition to $`\nu `$. This two-parameter family of temporally inhomogeneous processes indexed by $`(\nu ,\kappa ),\nu >1,\kappa [0,2(\nu +1))`$ is equivalent with the family of processes already studied by Yor. He called them the generalized meanders .
From the view-point of random matrix theory, studying time-development of stochastic systems by calculating, for example, the multitime correlation functions corresponds to considering multi-matrix models. In particular, the temporally inhomogeneous processes will be identified with such matrix models that matrices with different symmetries are coupled in a chain . Determination of all multitime correlation functions of systems, which allows us to determine scaling limits associated with the infinity limit of matrix sizes (i.e. the infinite-particle limit) is one of the main topics of the modern theory of RM . The finite and infinite particle systems showing the orthogonal-unitary and symplectic-unitary transitions, and transitions between class C to class CI were studied and multitime correlation functions were determined by Forrester, Nagao and Honner (FNH) , and by Nagao , respectively. The system in the Laguerre ensemble with $`\beta =1`$ initial condition reported in the former paper can be regarded as the $`\nu =\kappa \{0,1,2,\mathrm{}\}`$ case of the non-colliding system of the generalized meanders and the system reported in the latter paper as the $`(\nu ,\kappa )=(1/2,1)`$ case.
If we think about the system of generalized meanders apart from the RM theory, however, we can consider the parameters $`\nu `$ and $`\kappa `$ as real numbers, and not necessarily integers nor half-integers. In the present paper, we calculate the multitime correlation functions of non-colliding systems of (squared) generalized meanders for arbitrary values of parameters, provided they satisfy the condition $`\nu >1,\kappa [0,2(\nu +1))`$ so that the systems are not collapsed. We first define the $`N`$ particle systems in a finite time-interval $`[0,T]`$ and take the $`N=T\mathrm{}`$ limit to construct the two-parameter family of infinite particle systems. We prove that the multitime characteristic functions is given by a Fredholm Pfaffian and thus any multitime correlation function is given by a Pfaffian. Similarly to the results by FNH and Nagao and their temporally-homogeneous version (the determinantal process with the extended Bessel kernel ), the elements of the matrix kernels of Pfaffians are expressed using the Bessel functions, but we clarify the fact that they are generally given by the Riemann-Liouville differintegrals of the functions comprising the Bessel functions, which are used in fractional calculus (see, for example, ). This structure will explain the origin of the multiple integral expressions for the elements of the matrix kernels reported by FNH and Nagao .
The paper is organized as follows. In Section 2, the definitions of the generalized meanders of Yor and their non-colliding systems are given and the Riemann-Liouville differintegrals of the Bessel functions with appropriate factors are introduced. The main theorem for the infinite particle limit (Theorem 2.1) is then given. It is demonstrated that, if we take a further limit in the system of Theorem 2.1, we will obtain the temporally homogeneous system of infinite number of particles, which is a determinantal process with the extended Bessel kernel studied in (see also ). Using the properties of the Riemann-Liouville differintegrals, we show that Theorem 2.1 includes the results by FNH and Nagao as special cases. Section 3 is devoted to prove that for any finite number of particles $`N`$, the present system is a Pfaffian process (Theorem 3.1), in the sense that any multitime correlation function is given by a Pfaffian . These Pfaffian processes may be regarded as the continuous space-time version of the Pfaffian point processes and Pfaffian Schur processes studied by Borodin and Rains . Soshnikov used the term Pfaffian ensembles in . See also in the context of study of nonequilibrium phenomena in the polynuclear growth models, and in that of shape fluctuations of crystal facets. The processes studied in are also Pfaffian processes, since the ‘quaternion determinantal expressions’ of correlation functions, introduced and developed by Dyson, Mehta, Forrester, and Nagao , are readily transformed to Pfaffian expressions. The method of skew-orthogonal functions associated with the Laguerre polynomials are used in Section 4 in order to perform matrix inversion and give explicit expressions for the elements of matrix kernels of Pfaffians. Asymptotics in $`T=N\mathrm{}`$ are studied in Section 5. Appendices are given to show proofs of formulae and lemmas used in the text.
At the end of this introduction, we would like to refer to the papers , which reported the further extensions of RM theory in physics and the representation theory. We hope that the present paper will demonstrate the fruitfulness of developing the probability theory of interacting infinite particle systems in connection with the extensive study of (multi-)matrix models in the RM theory.
## 2 Definition of Processes and Results
### 2.1 Non-colliding systems of generalized meanders
Let $``$ and $``$ be the sets of integers and real numbers, respectively, and set $`=\{1,2,\mathrm{}\}`$, $`_0=\{0\}`$, $`_{}=_0`$, and $`_+=\{x:x0\}`$. Let $`\mathrm{\Gamma }(c),c(_{}\{0\})`$, be the Gamma function: $`\mathrm{\Gamma }(c)=_0^{\mathrm{}}𝑑ye^yy^{c1}`$ for $`c>0`$, and $`\mathrm{\Gamma }(c)=\mathrm{\Gamma }(c+[c]+1)/\{c(c+1)\mathrm{}(c+[c])\}`$ for $`c(\mathrm{},0)_{}`$, where $`[c]`$ is the largest integer that is less than or equal to the real number $`c`$. For $`t>0`$, $`x,y_+`$ and $`\nu >1`$ we denote by $`G_t^{(\nu )}(t;y|x)`$ the transition probability density of a $`2(\nu +1)`$-dimensional Bessel process ,
$`G^{(\nu )}(t;y|x)`$ $`=`$ $`{\displaystyle \frac{y^{\nu +1}}{x^\nu }}{\displaystyle \frac{1}{t}}e^{(x^2+y^2)/2t}I_\nu \left({\displaystyle \frac{xy}{t}}\right),x>0,y_+,`$
$`G^{(\nu )}(t;y|0)`$ $`=`$ $`{\displaystyle \frac{y^{2\nu +1}}{2^\nu \mathrm{\Gamma }(\nu +1)t^{\nu +1}}}e^{y^2/2t},y_+,`$
where $`I_\nu (z)`$ is the modified Bessel function : $`I_\nu (z)=_{n=0}^{\mathrm{}}(z/2)^{2n+\nu }/\{\mathrm{\Gamma }(n+1)\mathrm{\Gamma }(\nu +n+1)\}.`$ For $`T>0`$, $`\kappa [0,2(\nu +1))`$, we put
$$h_T^{(\nu ,\kappa )}(t,x)=_0^{\mathrm{}}𝑑yG^{(\nu )}(Tt;y|x)y^\kappa ,x_+,t[0,T],$$
and
$`G_T^{(\nu ,\kappa )}(s,x;t,y)={\displaystyle \frac{1}{h_T^{(\nu ,\kappa )}(s,x)}}G^{(\nu )}(ts;y|x)h_T^{(\nu ,\kappa )}(t,y),x>0,y_+,`$ (2.1)
$`G_T^{(\nu ,\kappa )}(0,0;t,y)={\displaystyle \frac{\mathrm{\Gamma }(\nu +1)}{\mathrm{\Gamma }(\nu +1\kappa /2)}}(2T)^{\kappa /2}G^{(\nu )}(t;y|0)h_T^{(\nu ,\kappa )}(t,y),y_+,`$ (2.2)
for $`0stT`$. This transition probability density $`G_T^{(\nu ,\kappa )}(s,x;t,y)`$ defines the temporally inhomogeneous process in a finite time-interval $`[0,T]`$, which is called a generalized meander. In particular, when $`\nu =1/2`$ and $`\kappa =1`$, it is identified with the process called a Brownian meander (see Chapter 3 in Yor ).
Now we consider the $`N`$-particle system of generalized meanders conditioned that they never collide in a finite time-interval $`[0,T]`$. Let
$$_{+<}^N=\{𝐱=(x_1,x_2,\mathrm{},x_N)_+^N:0x_1<x_2<\mathrm{}<x_N\}.$$
According to the determinantal formula of Karlin and McGregor , the transition probability density is given as
$$g_{N,T}^{(\nu ,\kappa )}(s,𝐱;t,𝐲)=\frac{f_{N,T}^{(\nu ,\kappa )}(s,𝐱;t,𝐲)𝒩_{N,T}^{(\nu ,\kappa )}(Tt,𝐲)}{𝒩_{N,T}^{(\nu ,\kappa )}(Ts,𝐱)},0stT,𝐱,𝐲_{+<}^N,$$
(2.3)
where
$$f_{N,T}^{(\nu ,\kappa )}(s,𝐱;t,𝐲)=\underset{1j,kN}{det}\left[G_T^{(\nu ,\kappa )}(s,x_j,t,y_k)\right],𝒩_{N,T}^{(\nu ,\kappa )}(t,𝐱)=_{_{+<}^N}𝑑𝐲f_{N,T}^{(\nu ,\kappa )}(Tt,𝐱;T,𝐲).$$
Since $`h_T^{(\nu ,0)}(t,x)=1`$, $`G_T^{(\nu ,0)}(s,x;t,y)=G^{(\nu )}(ts;y|x)`$ and thus $`f_{N,T}^{(\nu ,0)}`$ is temporally homogeneous and independent of $`T`$, we will write $`f_N^{(\nu )}(ts;𝐲|𝐱)`$ for $`f_{N,T}^{(\nu ,0)}(s,𝐱;t,𝐲)`$. Moreover, note that
$$f_{N,T}^{(\nu ,\kappa )}(s,𝐱;t,𝐲)=\frac{1}{h_T^{(\nu ,\kappa )}(s,𝐱)}f_N^{(\nu )}(ts;𝐲|𝐱)h_T^{(\nu ,\kappa )}(t,𝐲),$$
where $`h_T^{(\nu ,\kappa )}(t,𝐱)_{j=1}^Nh_T^{(\nu ,\kappa )}(t,x_j)`$ and $`h_T^{(\nu ,\kappa )}(T,𝐱)=_{j=1}^Nx_j^\kappa `$. Then (2.3) can be written as
$$g_{N,T}^{(\nu ,\kappa )}(s,𝐱;t,𝐲)=\frac{1}{\stackrel{~}{𝒩}_N^{(\nu ,\kappa )}(Ts,𝐱)}f_N^{(\nu )}(ts;𝐲|𝐱)\stackrel{~}{𝒩}_N^{(\nu ,\kappa )}(Tt,𝐲),$$
(2.4)
where
$$\stackrel{~}{𝒩}_N^{(\nu ,\kappa )}(t,𝐱)=_{_{+<}^N}𝑑𝐲f_N^{(\nu )}(t;𝐲|𝐱)\underset{j=1}{\overset{N}{}}y_j^\kappa .$$
(2.5)
In our previous paper it was shown that, taking the limit $`𝐱\mathrm{𝟎}(0,0,\mathrm{},0)`$ at the initial time $`s=0`$, (2.4) becomes
$`g_{N,T}^{(\nu ,\kappa )}(0,\mathrm{𝟎};t,𝐲)=C_{N,T}^{\nu ,\kappa }(t){\displaystyle \underset{j=1}{\overset{N}{}}}G^{(\nu )}(t,y_j|0){\displaystyle \underset{1j<kN}{}}(y_k^2y_j^2)\stackrel{~}{𝒩}_N^{(\nu ,\kappa )}(Tt,𝐲)`$
(2.6)
for $`\nu >1`$ and $`\kappa [0,2(\nu +1))`$, where
$$C_{N,T}^{\nu ,\kappa }(t)=\frac{T^{(N+\kappa 1)N/2}t^{(N1)N}}{2^{N(N\kappa 1)/2}}\underset{j=1}{\overset{N}{}}\frac{\mathrm{\Gamma }(\nu +1)\mathrm{\Gamma }(1/2)}{\mathrm{\Gamma }\left(j/2\right)\mathrm{\Gamma }\left((j+1+2\nu \kappa )/2\right)}.$$
The $`N`$-particle system of non-colliding generalized meanders all starting from the origin $`\mathrm{𝟎}`$ at time $`0`$ is defined by the transition probability density $`g_{N,T}^{(\nu ,\kappa )}`$ given above and it will be denoted by $`𝐗(t)_{+<}^N,t[0,T]`$ in the present paper. It makes a two-parameter family of temporally inhomogeneous processes parameterized by $`\nu >1`$ and $`\kappa [0,2(\nu +1))`$.
We denote by $`𝔛`$ the space of countable subsets $`\xi `$ of $``$ satisfying $`\mathrm{}(\xi K)<\mathrm{}`$ for any compact subset $`K`$. For $`𝐱=(x_1,x_2,\mathrm{},x_n)_{\mathrm{}=1}^{\mathrm{}}^{\mathrm{}}`$, we denote $`\{x_i\}_{i=1}^n𝔛`$ simply by $`\{𝐱\}`$. Then $`\mathrm{\Xi }_N^𝐗(t)=\{𝐗(t)\},t[0,T]`$, is the diffusion process on the set $`𝔛`$ with transition density function $`𝔤_{N,T}^{(\nu ,\kappa )}(s,\xi ;t,\eta )`$, $`0stT`$:
$$𝔤_{N,T}^{(\nu ,\kappa )}(s,\xi ;t,\eta )=\{\begin{array}{cc}g_{N,T}^{(\nu ,\kappa )}(s,𝐱;t,𝐲),\hfill & \text{if}s>0,\mathrm{}\xi =\mathrm{}\eta =N,\hfill \\ g_{N,T}^{(\nu ,\kappa )}(0,\mathrm{𝟎};t,𝐲),\hfill & \text{if}s=0,\xi =\{0\},\mathrm{}\eta =N,\hfill \\ 0,\hfill & \text{otherwise},\hfill \end{array}$$
where $`𝐱`$ and $`𝐲`$ are the elements of $`_{+<}^N`$ with $`\xi =\{𝐱\}`$, $`\eta =\{𝐲\}`$.
For the given time interval $`[0,T]`$, we consider the $`M`$ intermediate times $`0<t_1<t_2<\mathrm{}<t_M<T`$. For convenience, we set $`t_0=0`$, $`t_{M+1}=T`$. For $`𝐱^{(m)}^N`$, $`1mM+1`$, and $`N^{}=1,2,\mathrm{},N`$, we put $`𝐱_N^{}^{(m)}=(x_1^{(m)},x_2^{(m)},\mathrm{},x_N^{}^{(m)})`$ and $`\xi _m^N^{}=\{𝐱_N^{}^{(m)}\}`$. Then the multitime transition density function of the process $`\mathrm{\Xi }_N^𝐗(t)`$ is given by
$$𝔤_{N,T}^{(\nu ,\kappa )}(0,\{0\};t_1,\xi _1^N;\mathrm{};t_{M+1},\xi _{M+1}^N)=\underset{m=0}{\overset{M}{}}𝔤_{N,T}^{(\nu ,\kappa )}(t_m,\xi _m^N;t_{m+1},\xi _{m+1}^N),$$
(2.7)
where we assume $`\xi _0^N=\{0\}`$. For a sequence $`\{N_m\}_{m=1}^{M+1}`$ of positive integers less than or equal to $`N`$, we define the $`(N_1,N_2,\mathrm{},N_{M+1})`$-multitime correlation function by
$`\rho _{N,T}^𝐗(t_1,\xi _1^{N_1};t_2,\xi _2^{N_2};\mathrm{};t_{M+1},\xi _{M+1}^{N_{M+1}})`$
$`={\displaystyle \underset{_{m=1}^{M+1}_+^{NN_m}}{}}{\displaystyle \underset{m=1}{\overset{M+1}{}}}{\displaystyle \frac{1}{(NN_m)!}}{\displaystyle \underset{j=N_m+1}{\overset{N}{}}}dx_j^{(m)}𝔤_{N,T}^{(\nu ,\kappa )}(0,\{0\};t_1,\xi _1^N;\mathrm{};t_{M+1},\xi _{M+1}^N).`$ (2.8)
Associated with the generalized meander (2.1), (2.2), we consider a temporally inhomogeneous diffusion process with transition probability density
$`p_T^{(\nu ,\kappa )}(0,0;t,y)`$ $``$ $`G_T^{(\nu ,\kappa )}(0,0;t,\sqrt{y})\times {\displaystyle \frac{1}{2}}y^{1/2},y_+,`$
$`p_T^{(\nu ,\kappa )}(s,x;t,y)`$ $``$ $`G_T^{(\nu ,\kappa )}(s,\sqrt{x}:t,\sqrt{y})\times {\displaystyle \frac{1}{2}}y^{1/2},x>0,y_+,`$
$`t[0,T]`$, and call it a squared generalized meander. The $`N`$-particle system of non-colliding squared generalized meanders $`𝐘(t),t[0,T]`$, is then defined by
$$𝐘(t)=(X_1(t)^2,X_2(t)^2,\mathrm{},X_N(t)^2),t[0,T].$$
The correlation function $`\rho _{N,T}^𝐘`$ of $`\mathrm{\Xi }_N^𝐘(t)=\{𝐘(t)\}`$ is obtained from (2.8) through the relation
$`\rho _{N,T}^𝐘(t_1,\zeta _1^{N_1};t_2,\zeta _2^{N_2};\mathrm{};t_{M+1},\zeta _{M+1}^{N_{M+1}})`$
$`=\rho _{N,T}^𝐗(t_1,\xi _1^{N_1};t_2,\xi _2^{N_2};\mathrm{};t_{M+1},\xi _{M+1}^{N_{M+1}}){\displaystyle \underset{m=1}{\overset{M+1}{}}}{\displaystyle \underset{j=1}{\overset{N_m}{}}}{\displaystyle \frac{1}{2x_j^{(m)}}},`$ (2.9)
where $`\xi _m^{N_m}=\{𝐱_{N_m}^{(m)}\},\zeta _m^{N_m}=\{𝐲_{N_m}^{(m)}\}`$ with $`x_j^{(m)}=\sqrt{y_j^{(m)}}`$, $`1jN_m`$, $`1mM+1`$.
### 2.2 Riemann-Liouville differintegrals of Bessel functions
We consider the following left and right Riemann-Liouville differintegrals for integrable functions $`f`$ on $`_+`$,
$`{}_{0}{}^{}𝐃_{x}^{c}f(x)`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Gamma }(nc)}}\left({\displaystyle \frac{d}{dx}}\right)^n{\displaystyle _0^x}(xy)^{nc1}f(y)𝑑y,`$ (2.10)
$`{}_{x}{}^{}𝐃_{\mathrm{}}^{c}f(x)`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Gamma }(nc)}}\left({\displaystyle \frac{d}{dx}}\right)^n{\displaystyle _x^{\mathrm{}}}(yx)^{nc1}f(y)𝑑y,`$ (2.11)
where $`c`$ and $`n=[c+1]_+`$ with the notation $`x_+=\mathrm{max}\{x,0\}`$. It is easy to confirm that, if $`c_0`$, both of them are reduced to the ordinary multiple derivative,
$${}_{0}{}^{}𝐃_{x}^{c}f(x)=(1)_x^c𝐃_{\mathrm{}}^cf(x)=\left(\frac{d}{dx}\right)^cf(x),$$
and, if $`c_{}`$, they are equal to the multiple integrals,
$`{}_{0}{}^{}𝐃_{x}^{c}f(x)`$ $`=`$ $`{\displaystyle _0^x}𝑑y_{|c|1}{\displaystyle _0^{y_{|c|1}}}𝑑y_{|c|2}\mathrm{}{\displaystyle _0^{y_2}}𝑑y_1{\displaystyle _0^{y_1}}𝑑y_0f(y_0),`$
$`{}_{x}{}^{}𝐃_{\mathrm{}}^{c}f(x)`$ $`=`$ $`{\displaystyle _x^{\mathrm{}}}𝑑y_{|c|1}{\displaystyle _{y_{|c|1}}^{\mathrm{}}}𝑑y_{|c|2}\mathrm{}{\displaystyle _{y_2}^{\mathrm{}}}𝑑y_1{\displaystyle _{y_1}^{\mathrm{}}}𝑑y_0f(y_0).`$
For $`c(\mathrm{},0)_{}`$ (2.10) and (2.11) define fractional integrals, and for $`c_+_0`$ fractional differentials. The Riemann-Liouville differintegrals are most often used in the fractional calculus (see, for example, ).
Let $`J_\nu (z)`$ be the Bessel functions: $`J_\nu (z)=_{\mathrm{}=0}^{\mathrm{}}(1)^{\mathrm{}}(z/2)^{2\mathrm{}+\nu }/\{\mathrm{\Gamma }(\nu +\mathrm{}+1)\mathrm{}!\}.`$ We define functions $`\stackrel{~}{J}_\nu `$ and $`\widehat{J}_\nu `$ as
$`\stackrel{~}{J}_\nu (\theta ,\eta ,x,s)=(\theta \eta x)^{\nu /2}J_\nu (2\sqrt{\theta \eta x})e^{2s\theta \eta }=e^{2s\theta \eta }{\displaystyle \underset{\mathrm{}=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^{\mathrm{}}(\theta \eta x)^{\mathrm{}+\nu }}{\mathrm{\Gamma }(\nu +\mathrm{}+1)\mathrm{}!}},`$ (2.12)
$`\widehat{J}_\nu (\theta ,\eta ,x,s)=(\theta \eta x)^{\nu /2}J_\nu (2\sqrt{\theta \eta x})e^{2s\theta \eta }=e^{2s\theta \eta }{\displaystyle \underset{\mathrm{}=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(\theta \eta x)^{\mathrm{}}}{\mathrm{\Gamma }(\nu +\mathrm{}+1)\mathrm{}!}}.`$ (2.13)
We will use the following abbreviations for the Riemann-Liouville differintegrals of order $`c`$ of $`\stackrel{~}{J}_\nu `$ and $`\widehat{J}_\nu `$,
$`\stackrel{~}{J}_\nu ^{(c)}(\theta ,\eta ,x,s)`$ $`=`$ $`{}_{0}{}^{}𝐃_{\eta }^{c}\stackrel{~}{J}_\nu (\theta ,\eta ,x,s),\theta ,\eta >0,s,`$ (2.14)
$`\widehat{J}_\nu ^{(c)}(\theta ,\eta ,x,s)`$ $`=`$ $`{}_{\eta }{}^{}𝐃_{\mathrm{}}^{c}\widehat{J}_\nu (\theta ,\eta ,x,s),\theta ,\eta >0,s<0.`$ (2.15)
We note that, if $`c_0`$, $`\stackrel{~}{J}_\nu ^{(c)}`$ can be expanded as
$$\stackrel{~}{J}_\nu ^{(c)}(\theta ,\eta ,x,s)=\frac{1}{\mathrm{\Gamma }(c)}\underset{n=0}{\overset{\mathrm{}}{}}\frac{(1)^n\eta ^{nc}}{n!(nc)}\stackrel{~}{J}_\nu ^{(n)}(\theta ,\eta ,x,s),\theta ,\eta >0,s.$$
(2.16)
It is also noted that, since $`\widehat{J}_\nu (\theta ,\eta ,x,s)0`$ exponentially fast as $`\eta \mathrm{}`$, if $`s\theta <0`$,
$$\widehat{J}_\nu ^{(c)}(\theta ,\eta ,x,s)=\frac{1}{\mathrm{\Gamma }(nc)}_\eta ^{\mathrm{}}𝑑\xi (\xi \eta )^{nc1}\widehat{J}_\nu ^{(n)}(\theta ,\xi ,x,s),\theta ,\eta >0,s<0,$$
(2.17)
where $`n=[c+1]_+`$.
### 2.3 Results
We put
$$𝔞=𝔞(\nu ,\kappa )=\nu \frac{\kappa }{2},𝔟=𝔟(\nu ,\kappa )=\nu \kappa ,$$
(2.18)
and introduce functions $`𝒟(s,x;t,y)`$, $`\stackrel{~}{}(s,x;t,y)`$, $`𝒮(s,x;t,y)`$ and $`\stackrel{~}{𝒮}(s,x;t,y)`$, $`x,y_+`$, $`s,t<0`$,
$`𝒟(s,x;t,y)`$ $`=`$ $`{\displaystyle \frac{1}{4(xy)^{\kappa /2}}}{\displaystyle _0^1}d\theta \theta ^{1\kappa }[\stackrel{~}{J}_\nu ^{(𝔟1)}(\theta ,1,x,s)\stackrel{~}{J}_\nu ^{(𝔟)}(\theta ,1,y,t)`$
$`\stackrel{~}{J}_\nu ^{(𝔟)}(\theta ,1,x,s)\stackrel{~}{J}_\nu ^{(𝔟1)}(\theta ,1,y,t)],`$
$`\stackrel{~}{}(s,x;t,y)`$ $`=`$ $`(xy)^{\kappa /2}{\displaystyle _1^{\mathrm{}}}d\theta \theta ^{\kappa 1}[{\displaystyle _1^{\mathrm{}}}d\xi \xi ^𝔞\widehat{J}_\nu ^{(𝔟+1)}(\theta ,\xi ,x,s)\widehat{J}_\nu ^{(𝔟+1)}(\theta ,1,y,t)`$
$`\widehat{J}_\nu ^{(𝔟+1)}(\theta ,1,x,s){\displaystyle _1^{\mathrm{}}}d\xi \xi ^𝔞\widehat{J}_\nu ^{(𝔟+1)}(\theta ,\xi ,y,t)],`$
$`𝒮(s,x;t,y)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{x}{y}}\right)^{\kappa /2}{\displaystyle _0^1}d\theta [\widehat{J}_\nu ^{(𝔟+1)}(\theta ,1,x,s)\stackrel{~}{J}_\nu ^{(𝔟1)}(\theta ,1,y,t)`$ (2.19)
$`\{𝔞\stackrel{~}{J}_\nu ^{(𝔟1)}(\theta ,1,y,t)\stackrel{~}{J}_\nu ^{(𝔟)}(\theta ,1,y,t)\}{\displaystyle _1^{\mathrm{}}}d\xi \xi ^𝔞\widehat{J}_\nu ^{(𝔟+1)}(\theta ,\xi ,x,s)],`$
and
$`\stackrel{~}{𝒮}(s,x;t,y)=𝒮(s,x;t,y)\mathrm{𝟏}_{(s<t)}\left({\displaystyle \frac{y}{x}}\right)^{𝔟/2}𝒢(s,x;t,y),`$ (2.20)
where $`\mathrm{𝟏}_{(\omega )}`$ is the indicator function: $`\mathrm{𝟏}_{(\omega )}=1`$ if $`\omega `$ is satisfied and $`\mathrm{𝟏}_{(\omega )}=0`$ otherwise, and
$$𝒢(s,x;t,y)=_0^{\mathrm{}}𝑑\theta J_\nu (2\sqrt{\theta x})J_\nu (2\sqrt{\theta y})e^{2(st)\theta }.$$
(2.21)
For an integer $`N`$ and a skew-symmetric $`2N\times 2N`$ matrix $`A=(a_{ij})`$, the Pfaffian is defined as
$$\mathrm{Pf}(A)=\mathrm{Pf}_{1i<j2N}(a_{ij})=\frac{1}{N!}\underset{\sigma }{}\mathrm{sgn}(\sigma )a_{\sigma (1)\sigma (2)}a_{\sigma (3)\sigma (4)}\mathrm{}a_{\sigma (2N1)\sigma (2N)},$$
(2.22)
where the summation is extended over all permutations $`\sigma `$ of $`(1,2,\mathrm{},2N)`$ with restriction $`\sigma (2k1)<\sigma (2k),k=1,2,\mathrm{},N`$. We put
$$\widehat{\mathrm{\Xi }}_N^𝐘(s)=\{Y_1(T_N+s),Y_2(T_N+s),\mathrm{},Y_N(T_N+s)\},s[T_N,0),$$
and $`\widehat{\mathrm{\Xi }}_N^𝐘(s)=\{0\}`$, $`s(\mathrm{},T_N)`$. Then we can state the main theorem in the present paper.
###### Theorem 2.1
Let $`T_N=N`$. Then the process $`\widehat{\mathrm{\Xi }}_N^𝐘(s),s(\mathrm{},0)`$ converges to the process $`\widehat{\mathrm{\Xi }}_{\mathrm{}}^𝐘(s),s(\mathrm{},0)`$, as $`N\mathrm{}`$, in the sense of finite dimensional distributions, whose correlation functions $`\rho ^𝐘`$ are given by
$$\rho ^𝐘(s_1,\{𝐲_{N_1}^{(1)}\};s_2,\{𝐲_{N_2}^{(2)}\};\mathrm{};s_M,\{𝐲_{N_M}^{(M)}\})=\mathrm{Pf}\left[𝒜(𝐲_{N_1}^{(1)},𝐲_{N_2}^{(2)},\mathrm{},𝐲_{N_M}^{(M)})\right],$$
for any $`M1`$, any sequence $`\{N_m\}_{m=1}^M`$ of positive integers, and any strictly increasing sequence $`\{s_m\}_{m=1}^{M+1}`$ of nonpositive numbers with $`s_{M+1}=0`$, where $`𝒜(𝐲_{N_1}^{(1)},𝐲_{N_2}^{(2)},\mathrm{},𝐲_{N_M}^{(M)})`$ is the $`2_{m=1}^MN_m\times 2_{m=1}^MN_m`$ skew-symmetric matrix defined by
$$𝒜(𝐲_{N_1}^{(1)},𝐲_{N_2}^{(2)},\mathrm{},𝐲_{N_M}^{(M)})=\left(𝒜^{m,n}(y_i^{(m)},y_j^{(n)})\right)_{1iN_m,1jN_n,1m,nM}$$
with $`2\times 2`$ matrices $`𝒜^{m,n}(x,y)`$ ;
$`𝒜^{m,n}(x,y)=\left(\begin{array}{cc}𝒟(s_m,x;s_n,y)& \stackrel{~}{𝒮}(s_n,y;s_m,x)\\ \stackrel{~}{𝒮}(s_m,x;s_n,y)& \stackrel{~}{}(s_m,x;s_n,y)\end{array}\right).`$ (2.25)
In the infinite-particle system defined by Theorem 2.1, we can take the further limit:
$$s_m\mathrm{}\text{with the time differences }s_ns_m\text{ fixed},1m,nM.$$
In this limit, $`𝒟(s_m,x;s_n,y)\stackrel{~}{}(s_m,x;s_n,y)0`$, $`1m,nM`$, as we show in Appendix E. Therefore, we can replace $`𝒟`$ and $`\stackrel{~}{}`$ by zeros in the matrices. Then the Pfaffian is reduced to an ordinary determinant of the $`_{m=1}^{M+1}N_m\times _{m=1}^{M+1}N_m`$ matrix, $`𝔸(𝐲_{N_1}^{(1)},𝐲_{N_2}^{(2)},\mathrm{},𝐲_{N_M}^{(M)})=\left(𝐚^{m,n}(y_i^{(m)},y_j^{(n)})\right)_{1iN_m,1jN_n,1m,nM}`$ with the elements
$$𝐚^{m,n}(y_i^{(m)},y_j^{(n)})=\stackrel{~}{𝕊}(s_m,y_i^{(m)};s_n,y_j^{(n)}),$$
where
$`\stackrel{~}{𝕊}(s,x;t,y)=\{\begin{array}{cc}{\displaystyle _0^1}𝑑\theta J_\nu (2\sqrt{\theta x})J_\nu (2\sqrt{\theta y})e^{2(st)\theta },\hfill & \text{if}s>t,\hfill \\ & \\ {\displaystyle \frac{J_\nu (2\sqrt{x})\sqrt{y}J_\nu ^{}(2\sqrt{y})J_\nu (2\sqrt{y})\sqrt{x}J_\nu ^{}(2\sqrt{x})}{xy}},\hfill & \text{if}s=t,\hfill \\ & \\ {\displaystyle _1^{\mathrm{}}}𝑑\theta J_\nu (2\sqrt{\theta x})J_\nu (2\sqrt{\theta y})e^{2(st)\theta },\hfill & \text{if}s<t,\hfill \end{array}`$ (2.31)
with $`J_\nu ^{}=dJ_\nu (z)/dz`$. Hence, in this limit we obtain a temporally homogeneous system of infinite number of particles, whose correlation functions are given by
$$\stackrel{~}{\rho }^𝐘(s_1,\{𝐲_{N_1}^{(1)}\};s_2,\{𝐲_{N_2}^{(2)}\};\mathrm{};s_M,\{𝐲_{N_M}^{(M)}\})=det𝔸(𝐲_{N_1}^{(1)},𝐲_{N_2}^{(2)},\mathrm{},𝐲_{N_M}^{(M)}).$$
(2.32)
Remark 1. Forrester, Nagao, and Honner studied the orthogonal-unitary and symplectic-unitary universality transitions in random matrix theory by giving the quaternion determinantal expressions of (two-time) correlation functions for parametric RM models. One of their results for the ‘Laguerre ensemble with $`\beta =1`$ initial condition’, which shows the orthogonal-unitary transition, can be reproduced from Theorem 2.1 by setting
$$(\mathrm{i})\kappa =\nu 𝔞=\frac{\nu }{2},𝔟=0,\text{where}\nu _0.$$
This fact may be readily seen, if we notice that by definition
$`\stackrel{~}{J}_\nu ^{(1)}(\theta ,1,x,s)`$ $`=`$ $`{\displaystyle _0^1}𝑑\eta (\theta \eta x)^{\nu /2}J_\nu (2\sqrt{\theta \eta x})e^{2s\theta \eta }`$
$`=`$ $`\theta ^1x^{\nu /2}{\displaystyle _0^\theta }𝑑uu^{\nu /2}J_\nu (2\sqrt{ux})e^{2su}.`$
Remark 2. Nagao’s result on the multitime correlation functions for vicious random walk with a wall can be regarded as the special case of Theorem 2.1, in which
$$(\mathrm{ii})\nu =\frac{1}{2},\kappa =1𝔞=0,𝔟=\frac{1}{2}.$$
This fact can be confirmed by noting that, by definition (2.14) with $`\stackrel{~}{J}_{1/2}(\theta ,0;x,s)=0`$,
$`\stackrel{~}{J}_{1/2}^{(1/2)}(\theta ,1,x,s)`$ $`=`$ $`{\displaystyle \frac{(\theta x)^{1/4}}{\sqrt{\pi }}}{\displaystyle _0^1}𝑑\eta (1\eta )^{1/2}\eta ^{1/4}J_{1/2}(2\sqrt{\theta \eta x})e^{2s\theta \eta },`$
$`\stackrel{~}{J}_{1/2}^{(1/2)}(\theta ,1,x,s)`$ $`=`$ $`{\displaystyle \frac{(\theta x)^{1/4}}{\sqrt{\pi }}}{\displaystyle _0^1}𝑑\eta (1\eta )^{1/2}{\displaystyle \frac{d}{d\eta }}\left\{\eta ^{1/4}J_{1/2}(2\sqrt{\theta \eta x})e^{2s\theta \eta }\right\},`$
by (2.17),
$$\widehat{J}_{1/2}^{(1/2)}(\theta ,\eta ,x,s)=\frac{(\theta x)^{1/4}}{\sqrt{\pi }}_\eta ^{\mathrm{}}𝑑\xi (\xi \eta )^{1/2}\frac{d}{d\xi }\left\{\xi ^{1/4}J_{1/2}(2\sqrt{\theta \xi x})e^{2s\theta \xi }\right\},s<0,$$
and, by definition (2.15),
$$_1^{\mathrm{}}𝑑\xi \widehat{J}_{1/2}^{(1/2)}(\theta ,\xi ,x,s)=\frac{(\theta x)^{1/4}}{\sqrt{\pi }}_1^{\mathrm{}}𝑑\eta (\eta 1)^{1/2}\eta ^{1/4}J_{1/2}(2\sqrt{\theta \eta x})e^{2s\theta \eta },s<0.$$
In this case, the system shows the transition between the class C and class CI of the Bogoliubov-de Gennes universality classes of nonstandard RM theory .
Remark 3. From the results for finite non-colliding processes , we expect that, when
$$(\mathrm{iii})\kappa =\nu +1𝔞=\frac{\nu 1}{2},𝔟=1,\text{where}\nu _0,$$
the present infinite particle system will show the transition from the chiral GUE to the chiral GOE of the universality classes and when
$$(\mathrm{iv})\nu =\frac{1}{2},\kappa =0𝔞=𝔟=\frac{1}{2},$$
that from the class D to the ‘real-component version’ of class D of the Bogoliubov-de Gennes universality classes .
Remark 4. Following the argument given in , tightness in time can be proved and transition phenomena observed in the limit $`s_M0`$ may be generally discussed, which will be reported elsewhere.
Remark 5. The homogeneous system (2.32) was studied in .
## 3 Correlation Functions Given by Pfaffians
### 3.1 The multitime transition density
If we put
$`\stackrel{~}{G}^{(\nu ,\kappa )}(t,y|x)=G^{(\nu )}(t,y|x)\times \left({\displaystyle \frac{y}{x}}\right)^\kappa ,x>0,y_+,`$
$`\stackrel{~}{G}^{(\nu ,\kappa )}(t,y|0)=G^{(\nu )}(t,y|0)\times y^\kappa ,y_+,`$
the multitime transition density (2.7) with $`t_0=0,t_{M+1}=T`$, and $`\xi _0=\{0\}`$ is written as
$`𝔤_{N,T}^{(\nu ,\kappa )}(0,\{0\};t_1,\{𝐱_N^{(1)}\};\mathrm{};t_{M+1},\{𝐱_N^{(M+1)}\})`$
$`=C_{N,T}^{\nu ,\kappa }(t_1){\displaystyle \underset{1j<kN}{}}\left\{(x_k^{(1)})^2(x_j^{(1)})^2\right\}{\displaystyle \underset{1j<kN}{}}\mathrm{sgn}(x_k^{(M+1)}x_j^{(M+1)})`$
$`\times {\displaystyle \underset{j=1}{\overset{N}{}}}\stackrel{~}{G}^{(\nu ,\kappa )}(t_1,x_j^{(1)}|0){\displaystyle \underset{m=1}{\overset{M}{}}}\underset{1j,kN}{det}\left[\stackrel{~}{G}^{(\nu ,\kappa )}(t_{m+1}t_m,x_j^{(m+1)}|x_k^{(m)})\right],`$
where (2.4) and (2.6) with (2.5) are used.
Through the relation (2.9), the multitime transition density for the process $`\{𝐘(t)\},t[0,T]`$, denoted by $`𝔭_{N,T}^{(\nu ,\kappa )}`$ is then written as
$`𝔭_{N,T}^{(\nu ,\kappa )}(0,\{0\};t_1,\{𝐲^{(1)}\};\mathrm{};t_{M+1},\{𝐲^{(M+1)}\})`$
$`=C_{N,T}^{\nu ,\kappa }(t_1)h_N(𝐲^{(1)})\mathrm{sgn}\left(h_N(𝐲^{(M+1)})\right)`$
$`\times {\displaystyle \underset{k=1}{\overset{N}{}}}\stackrel{~}{p}^{(\nu ,\kappa )}(t_1,y_k^{(1)}|0){\displaystyle \underset{m=1}{\overset{M}{}}}\underset{1j,kN}{det}\left[\stackrel{~}{p}^{(\nu ,\kappa )}(t_{m+1}t_m,y_j^{(m+1)}|y_k^{(m)})\right],`$ (3.1)
where
$$h_N(𝐲)\underset{1i<jN}{}(y_jy_i),𝐲^N,$$
$`\stackrel{~}{p}^{(\nu ,\kappa )}(t,y|0)`$ $``$ $`\stackrel{~}{G}^{(\nu ,\kappa )}(t,\sqrt{y}|0)\times {\displaystyle \frac{1}{2}}y^{1/2}`$
$`=`$ $`{\displaystyle \frac{y^𝔞}{2^{\nu +1}\mathrm{\Gamma }(\nu +1)t^{\nu +1}}}e^{y/2t},y_+,`$
$`\stackrel{~}{p}^{(\nu ,\kappa )}(ts,y|x)`$ $``$ $`\stackrel{~}{G}^{(\nu ,\kappa )}(ts,\sqrt{y}|\sqrt{x})\times {\displaystyle \frac{1}{2}}y^{1/2}`$ (3.2)
$`=`$ $`{\displaystyle \frac{e^{(x+y)/\{2(ts)\}}}{2(ts)}}\left({\displaystyle \frac{y}{x}}\right)^{𝔟/2}I_\nu \left({\displaystyle \frac{\sqrt{xy}}{ts}}\right),x>0,y_+.`$
Expectations related to the process $`\{𝐘(t_1)\},\{𝐘(t_2)\},\mathrm{},\{𝐘(t_{M+1})\}`$ are denote by $`𝔼_{N,T}^𝐘`$ :
$`𝔼_{N,T}^𝐘\left[f(\{𝐘(t_1)\},\{𝐘(t_2)\},\mathrm{},\{𝐘(t_{M+1})\})\right]=\left({\displaystyle \frac{1}{N!}}\right)^{M+1}{\displaystyle _{_+^{N(M+1)}}}{\displaystyle \underset{m=1}{\overset{M+1}{}}}d𝐲^{(m)}`$
$`\times f(\{𝐲^{(1)}\},\{𝐲^{(2)}\},\mathrm{},\{𝐲^{(M+1)}\})𝔭_{N,T}^{(\nu ,\kappa )}(0,\{0\};t_1,\{𝐲^{(1)}\};\mathrm{};t_{M+1},\{𝐲^{(M+1)}\}).`$ (3.3)
### 3.2 Fredholm Pfaffian representation of characteristic function and Pfaffian process
For simplicity of expressions, we assume from now on that the number of particles $`N`$ is even. The references will be useful to give necessary modifications to the following expressions in the case that $`N`$ is odd. Let $`C_0()`$ be the set of all continuous real functions with compact supports. For $`𝐟=(f_1,f_2,\mathrm{},f_{M+1})C_0()^{M+1}`$, and $`𝜽=(\theta _1,\theta _2,\mathrm{},\theta _{M+1})^{M+1}`$, the multitime characteristic function is defined for the process $`\{𝐘(t)\},t[0,T]`$ as
$$\mathrm{\Psi }_{N,T}^𝐘(𝐟;𝜽)=𝔼_{N,T}^𝐘\left[\mathrm{exp}\left\{\sqrt{1}\underset{m=1}{\overset{M+1}{}}\theta _m\underset{i_m=1}{\overset{N}{}}f_m(Y_{i_m}(t_m))\right\}\right]$$
(3.4)
Let $`\chi _m(x)=e^{\sqrt{1}\theta _mf_m(x)}1,\mathrm{\hspace{0.17em}1}mM+1.`$ Then by the definition of multitime correlation function (2.9) with (2.8), we have
$`\mathrm{\Psi }_{N,T}^𝐘(𝐟;𝜽)`$ $`=`$ $`{\displaystyle \underset{N_1=0}{\overset{N}{}}}{\displaystyle \underset{N_2=0}{\overset{N}{}}}\mathrm{}{\displaystyle \underset{N_{M+1}=0}{\overset{N}{}}}{\displaystyle \underset{m=1}{\overset{M+1}{}}}{\displaystyle \frac{1}{N_m!}}{\displaystyle _{_+^{N_1}}}𝑑𝐲_{N_1}^{(1)}{\displaystyle _{_+^{N_2}}}𝑑𝐲_{N_2}^{(2)}\mathrm{}{\displaystyle _{_+^{N_{M+1}}}}𝑑𝐲_{N_{M+1}}^{(M+1)}`$ (3.5)
$`\times {\displaystyle \underset{m=1}{\overset{M+1}{}}}{\displaystyle \underset{i^{(m)}=1}{\overset{N_m}{}}}\chi _m\left(y_{i^{(m)}}^{(m)}\right)\rho _{N,T}^𝐘(t_1,\{𝐲_{N_1}^{(1)}\};t_2,\{𝐲_{N_2}^{(2)}\};\mathrm{};t_{M+1},\{𝐲_{N_{M+1}}^{(M+1)}\}),`$
that is, the multitime characteristic function is a generating function of multitime correlation functions $`\rho _{N,T}^𝐘`$.
We consider a vector space $`𝒱`$ with the orthonormal basis $`\left\{|m,x\right\}_{1mM+1,x_+}`$, which satisfies
$$m,x|n,y=\delta _{mn}\delta (xy),m,n=1,2,\mathrm{},M+1,x,y_+,$$
(3.6)
where $`\delta _{mn}`$ and $`\delta (xy)`$ denote Kronecker’s delta and Dirac’s $`\delta `$-measure, respectively. We introduce the operators $`\widehat{J},\widehat{p},\widehat{p}_+,\widehat{p}_{}`$ and $`\widehat{\chi }`$ acting on $`𝒱`$ as follows
$`m,x|\widehat{J}|n,y`$ $`=`$ $`\mathrm{𝟏}_{(m=n=M+1)}\mathrm{sgn}(yx),`$ (3.7)
$`m,x|\widehat{p}|n,y`$ $`=`$ $`\mathrm{𝟏}_{(m<n)}\stackrel{~}{p}^{(\nu ,\kappa )}(t_nt_m,y|x)+\mathrm{𝟏}_{(m>n)}\stackrel{~}{p}^{(\nu ,\kappa )}(t_mt_n,x|y)`$ (3.8)
$`+`$ $`\mathrm{𝟏}_{(m=n)}\delta (xy),`$
$`m,x|\widehat{p}_+|n,y`$ $`=`$ $`\mathrm{𝟏}_{(m<n)}\stackrel{~}{p}^{(\nu ,\kappa )}(t_nt_m,y|x)=n,y|\widehat{p}_{}|m,x,`$ (3.9)
$`m,x|\widehat{\chi }|n,y`$ $`=`$ $`\chi _m(x)\delta _{mn}\delta (xy),`$ (3.10)
and we will use the convention
$$m,x|\widehat{A}|n,yn,y|\widehat{B}|\mathrm{},z)=\underset{n=1}{\overset{M+1}{}}__+𝑑yA(m,x;n,y)B(n,y;\mathrm{},z)=m,x|\widehat{A}\widehat{B}|\mathrm{},z$$
for operators $`\widehat{A}`$ and $`\widehat{B}`$ with $`m,x|\widehat{A}|n,y=A(m,x;n,y)`$ and $`m,x|\widehat{B}|n,y=B(m,x;n,y)`$.
Let $`M_i(x)`$ be an arbitrary polynomial of $`x`$ with degree $`i`$ in the form $`M_i(x)=b_ix^i+\mathrm{}`$ with a constant $`b_i0`$ for $`i_0`$. Since the product of differences $`h_N(𝐱)`$ is equal to the Vandermonde determinant, we have
$$h_N(𝐱)=\left\{\underset{k=1}{\overset{N}{}}b_{k1}\right\}^1\underset{1i,jN}{det}\left[M_{i1}(x_j)\right].$$
(3.11)
Then we consider the set of linearly independent vectors $`\left\{|i;i\right\}`$ in $`𝒱`$ defined by
$$|i=|m,xm,x|i,$$
where
$$m,x|i=i|m,x=__+𝑑yM_{i1}(y)\stackrel{~}{p}^{(\nu ,\kappa )}(t_1,y|0)\stackrel{~}{p}^{(\nu ,\kappa )}(t_mt_1,x|y),$$
(3.12)
$`i,m=1,2,\mathrm{},M+1,x_+`$. We will use the convention
$$i|\widehat{A}|jj|\widehat{B}|m,x=\underset{j=1}{\overset{\mathrm{}}{}}A_{ij}B_j^{(m)}(x)=i|\widehat{A}\widehat{B}|m,x,$$
for $`A_{ij}=i|\widehat{A}|j`$ and $`B_j^{(m)}(x)=j|\widehat{B}|m,x.`$ It should be noted that the vecors $`\left\{|i;i\right\}`$ are not assumed to be mutually orthogonal. By these vectors, however, any operator $`\widehat{A}`$ on $`𝒱`$ may have a semi-infinite matrix representation $`A=\left(i|\widehat{A}|j\right)_{i,j}`$. If the matrix $`A`$ representing an operator $`\widehat{A}`$ is invertible, we define the operator $`\widehat{A}^{}`$ so that its matrix representation is the inverse of $`A`$;
$$\left(i|\widehat{A}^{}|j\right)_{i,j}=A^1,$$
(3.13)
that is, $`i|\widehat{A}|jj|\widehat{A}^{}|k=i|\widehat{A}\widehat{A}^{}|k=\delta _{ik},i,k.`$
Let $`𝒫_N`$ be a linear operator projecting $`\mathrm{𝑆𝑝𝑎𝑛}\left\{|i;i\right\}`$ to its $`N`$-dimensional subspace $`\mathrm{𝑆𝑝𝑎𝑛}\{|i;i=1,2,\mathrm{},N\}`$ such that
$$i|𝒫_N|m,x=m,x|𝒫_N|i=\{\begin{array}{cc}i|m,x,\hfill & \text{if}1iN,\hfill \\ & \\ 0,\hfill & \text{otherwise}.\hfill \end{array}$$
We will use the abbreviation $`\widehat{A}_N=𝒫_N\widehat{A}𝒫_N`$ for an operator $`\widehat{A}`$. If the $`N\times N`$ matrix defined by $`A_N=(i|\widehat{A}_N|j)_{1i,jN}`$ is invertible, then $`(\widehat{A}_N)^{}`$ is defined so that $`\left(i|(\widehat{A}_N)^{}|j\right)_{1i,jN}=(A_N)^1`$, and $`i|(\widehat{A}_N)^{}|j=0`$, if $`iN+1`$ or $`jN+1`$.
As shown in A, we can prove that
$$\left\{\mathrm{\Psi }_{N,T}^𝐘(𝐟;𝜽)\right\}^2=\mathrm{Det}\left(I_2\delta _{mn}\delta (xy)+\left(\begin{array}{cc}\stackrel{~}{S}^{m,n}(x,y)& \stackrel{~}{I}^{m,n}(x,y)\\ D^{m,n}(x,y)& \stackrel{~}{S}^{n,m}(y,x)\end{array}\right)\chi _n(y)\right),$$
(3.14)
where $`\mathrm{Det}`$ denotes the Fredholm determinant. Here $`I_2`$ is the unit matrix with size 2,
$`D^{m,n}(x,y)`$ $`=`$ $`m,x|(\widehat{J}_N)^{}|n,y,`$
$`S^{m,n}(x,y)`$ $`=`$ $`m,x|\widehat{p}\widehat{J}(\widehat{J}_N)^{}|n,y,`$
$`I^{m,n}(x,y)`$ $`=`$ $`m,x|\widehat{p}\widehat{J}(\widehat{J}_N)^{}\widehat{J}\widehat{p}|n,y,`$ (3.15)
and
$`\stackrel{~}{S}^{m,n}(x,y)`$ $`=`$ $`S^{m,n}(x,y)m,x|\widehat{p}_+|n,y`$
$`\stackrel{~}{I}^{m,n}(x,y)`$ $`=`$ $`I^{m,n}(x,y)+m,x|\widehat{p}\widehat{J}\widehat{p}|n,y.`$ (3.16)
It implies that the multitime characteristic function is given by the Fredholm Pfaffian ,
$$\mathrm{\Psi }_{N,T}^𝐘(𝐟;𝜽)=\mathrm{PF}\left(J_2\delta _{mn}\delta (xy)+\sqrt{\chi _m(x)}A^{m,n}(x,y)\sqrt{\chi _n(y)}\right),$$
(3.17)
where $`J_2=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$ and
$`A^{m,n}(x,y)`$ $`=`$ $`J_2\left(\begin{array}{cc}\stackrel{~}{S}^{m,n}(x,y)& \stackrel{~}{I}^{m,n}(x,y)\\ D^{m,n}(x,y)& \stackrel{~}{S}^{n,m}(y,x)\end{array}\right)`$ (3.20)
$`=`$ $`\left(\begin{array}{cc}D^{m,n}(x,y)& \stackrel{~}{S}^{n,m}(y,x)\\ \stackrel{~}{S}^{m,n}(x,y)& \stackrel{~}{I}^{m,n}(x,y)\end{array}\right).`$ (3.23)
It is defined by
$`\mathrm{PF}\left(J_2\delta _{mn}\delta (xy)+\sqrt{\chi _m(x)}A^{m,n}(x,y)\sqrt{\chi _n(y)}\right)`$
$`={\displaystyle \underset{N_1=0}{\overset{N}{}}}{\displaystyle \underset{N_2=0}{\overset{N}{}}}\mathrm{}{\displaystyle \underset{N_{M+1}=0}{\overset{N}{}}}{\displaystyle \underset{m=1}{\overset{M+1}{}}}{\displaystyle \frac{1}{N_m!}}{\displaystyle _{_+^{N_1}}}𝑑𝐲_{N_1}^{(1)}{\displaystyle _{_+^{N_2}}}𝑑𝐲_{N_2}^{(2)}\mathrm{}{\displaystyle _{_+^{N_{M+1}}}}𝑑𝐲_{N_{M+1}}^{(M+1)}`$
$`{\displaystyle \underset{m=1}{\overset{M+1}{}}}{\displaystyle \underset{i^{(m)}=1}{\overset{N_m}{}}}\chi _m\left(y_{i^{(m)}}^{(m)}\right)\mathrm{Pf}\left(A(𝐲_{N_1}^{(1)},𝐲_{N_2}^{(2)},\mathrm{},𝐲_{N_{M+1}}^{(M+1)})\right),`$ (3.24)
where $`A(𝐲_{N_1}^{(1)},𝐲_{N_2}^{(2)},\mathrm{},𝐲_{N_{M+1}}^{(M+1)})`$ denotes the $`2_{m=1}^{M+1}N_m\times 2_{m=1}^{M+1}N_m`$ skew-symmetric matrices constructed from (3.23) as
$$A(𝐲_{N_1}^{(1)},𝐲_{N_2}^{(2)},\mathrm{},𝐲_{N_{M+1}}^{(M+1)})=\left(A^{m,n}(y_i^{(m)},y_j^{(n)})\right)_{1iN_m,1jN_n,1m,nM+1}$$
for $`N_m=1,2,\mathrm{},N,1mM+1`$. Comparison of (3.5) and (3.17) with (3.24) immediately gives the following statement.
###### Theorem 3.1
The $`N`$-particle non-colliding system of squared generalized meanders $`𝐘(t),t[0,T]`$ is a Pfaffian process, in the sense that any multitime correlation function is given by a Pfaffian
$$\rho _{N,T}^𝐘(t_1,\{𝐲_{N_1}^{(1)}\};t_2,\{𝐲_{N_2}^{(2)}\};\mathrm{};t_{M+1},\{𝐲_{N_{M+1}}^{(M+1)}\})=\mathrm{Pf}\left(A(𝐲_{N_1}^{(1)},𝐲_{N_2}^{(2)},\mathrm{},𝐲_{N_{M+1}}^{(M+1)})\right).$$
## 4 Skew-Orthogonal Functions and Matrix Inversion
### 4.1 Skew-symmetric inner products
Consider the $`N\times N`$ skew-symmetric matrix $`A_0=((A_0)_{ij})_{1i,jN}`$ with
$$(A_0)_{ij}=i|\widehat{J}_N|j=i|m,xm,x|\widehat{J}|n,yn,y|j,i,j=1,2,\mathrm{},N.$$
(4.1)
In order to clarify the fact that each element $`(A_0)_{ij}`$ is a functional of the polynomials $`M_{i1}(x)`$ and $`M_{j1}(x)`$ through (3.12), we introduce the skew-symmetric inner product
$$f,g_0^{\mathrm{}}𝑑x_0^{\mathrm{}}𝑑yF(x,y)\stackrel{~}{p}^{(\nu ,\kappa )}(t_1,x|0)\stackrel{~}{p}^{(\nu ,\kappa )}(t_1,y|0)f(x)g(y),$$
(4.2)
where
$$F(x,y)=_0^{\mathrm{}}𝑑w_0^w𝑑z\left|\begin{array}{cc}\stackrel{~}{p}^{(\nu ,\kappa )}(Tt_1,z|x)& \stackrel{~}{p}^{(\nu ,\kappa )}(Tt_1,w|x)\\ \stackrel{~}{p}^{(\nu ,\kappa )}(Tt_1,z|y)& \stackrel{~}{p}^{(\nu ,\kappa )}(Tt_1,w|y)\end{array}\right|,x,y_+.$$
(4.3)
Then we have the expression
$$(A_0)_{ij}=M_{i1},M_{j1},i,j=1,2,\mathrm{},N.$$
(4.4)
We now rewrite the skew-symmetric inner product (4.2) by using the simpler one
$`f,g_{}`$ $`=`$ $`g,f_{}`$ (4.5)
$``$ $`{\displaystyle _0^{\mathrm{}}}𝑑we^{w/2}w^𝔞{\displaystyle _0^w}𝑑ze^{z/2}z^𝔞\left\{f(z)g(w)f(w)g(z)\right\},`$
which we call the elementary skew-symmetric inner product. Remind that $`\stackrel{~}{p}^{(\nu ,\kappa )}`$ is given by (3.2) using the modified Bessel function. We will expand it in terms of the Laguerre polynomials, $`L_j^\alpha (x)=(x^\alpha e^x/j!)(d/dx)^j(e^xx^{j+\alpha })`$, $`\alpha `$, $`j_0`$, using the formula
$$\underset{j=0}{\overset{\mathrm{}}{}}\frac{\mathrm{\Gamma }(j+1)L_j^\nu (x)L_j^\nu (y)r^j}{\mathrm{\Gamma }(j+1+\nu )}=\frac{1}{1r}e^{\frac{(x+y)r}{1r}}(xyr)^{\nu /2}I_\nu \left(\frac{2\sqrt{xyr}}{1r}\right),|r|<1,\nu >1.$$
(4.6)
(See the corresponding calculation for the non-colliding Brownian particles in , where the heat kernel was expanded in terms of the Hermite polynomials.) For this purpose, it is useful to introduce the variables
$$c_n=\frac{t_n(2Tt_n)}{T},\chi _n=\frac{2Tt_n}{t_n},n=1,2,\mathrm{},M+1,$$
since we can see that
$`\stackrel{~}{p}^{(\nu ,\kappa )}(t_nt_m,c_n\eta |c_m\xi )={\displaystyle \frac{1}{2(t_nt_m)}}I_\nu \left({\displaystyle \frac{2\sqrt{\xi \eta \chi _n/\chi _m}}{1\chi _n/\chi _m}}\right)\left({\displaystyle \frac{c_n\eta }{c_m\xi }}\right)^{𝔟/2}`$
$`\times \mathrm{exp}\left[\left({\displaystyle \frac{1}{1\chi _n/\chi _m}}1+{\displaystyle \frac{t_m}{2T}}\right)\xi \left({\displaystyle \frac{1}{1\chi _n/\chi _m}}{\displaystyle \frac{t_n}{2T}}\right)\eta \right],`$
and, if we apply the formula (4.6) with $`r=\chi _n/\chi _m,x=\xi `$ and $`y=\eta `$, it is written as
$`\stackrel{~}{p}^{(\nu ,\kappa )}(t_nt_m,c_n\eta |c_m\xi )`$ $`=`$ $`\left({\displaystyle \frac{t_m}{t_n}}\right)^{\nu +1}c_m^{𝔞1}\xi ^{\kappa /2}(c_n\eta )^𝔞\mathrm{exp}\left[{\displaystyle \frac{t_m}{2T}}\xi (1{\displaystyle \frac{t_n}{2T}})\eta \right]`$ (4.7)
$`\times {\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{\Gamma }(j+1)}{\mathrm{\Gamma }(j+1+\nu )}}\left({\displaystyle \frac{\chi _n}{\chi _m}}\right)^jL_j^\nu (\xi )L_j^\nu (\eta ).`$
That is, $`c_n`$ and $`\chi _n`$ give the spatial scale of spread of $`N`$ particles and the proper temporal factor at time $`t_n`$, respectively. (See equation (17) and explanation below it in , where the variable $`c_n`$ was determined by showing that the one-particle density obeys Wigner’s semicircle law scaled by $`c_n`$ for the non-colliding Brownian particles.) In particular, for $`n=M+1`$ we have
$`\stackrel{~}{p}^{(\nu ,\kappa )}(Tt_m,T\eta |c_m\xi )`$ $`=`$ $`{\displaystyle \frac{t_m^{\nu +1}}{T^{\kappa /2+1}}}c_m^{𝔞1}\xi ^{\kappa /2}\eta ^𝔞\mathrm{exp}\left[\left(1{\displaystyle \frac{t_m}{T}}\right){\displaystyle \frac{\xi }{2}}\right]`$ (4.8)
$`\times e^{\xi /2}e^{\eta /2}{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{\Gamma }(j+1)}{\mathrm{\Gamma }(j+1+\nu )}}\chi _m^jL_j^\nu (\xi )L_j^\nu (\eta ),`$
since $`c_{M+1}=T`$ and $`\chi _{M+1}=1`$. Then we obtain the relation
$`f\left({\displaystyle \frac{}{c_1}}\right),g\left({\displaystyle \frac{}{c_1}}\right)`$ $`=`$ $`{\displaystyle \frac{2^{2\nu 2}T^\kappa }{\mathrm{\Gamma }(\nu +1)^2}}{\displaystyle _0^{\mathrm{}}}𝑑x{\displaystyle _0^{\mathrm{}}}𝑑ye^xe^yx^\nu y^\nu f(x)g(y)`$ (4.9)
$`\times `$ $`{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}\chi _1^{jk}L_j^\nu (x)L_k^\nu (y){\displaystyle \frac{\mathrm{\Gamma }(j+1)}{\mathrm{\Gamma }(j+1+\nu )}}L_j^\nu ,{\displaystyle \frac{\mathrm{\Gamma }(k+1)}{\mathrm{\Gamma }(k+1+\nu )}}L_k^\nu _{}.`$
### 4.2 Skew-orthogonal polynomials
For $`\alpha `$ and $`n`$ we define
$$\left(\genfrac{}{}{0pt}{}{n+\alpha }{n}\right)=\{\begin{array}{cc}\frac{\mathrm{\Gamma }(n+\alpha +1)}{\mathrm{\Gamma }(n+1)\mathrm{\Gamma }(\alpha +1)},\hfill & \text{if }n\text{}\alpha _{},\hfill \\ \frac{(1)^n\mathrm{\Gamma }(\alpha )}{\mathrm{\Gamma }(n+1)\mathrm{\Gamma }(n\alpha )},\hfill & \text{if }n\text{}n+\alpha _{},\hfill \\ 0,\hfill & \text{if }n\text{}\alpha _{}\text{}n+\alpha _0,\hfill \\ 1,\hfill & \text{if }n=0,\hfill \\ 0,\hfill & \text{if }n_{}.\hfill \end{array}$$
(4.10)
Note that for $`n,\alpha _{}`$ with $`n+\alpha 1`$, $`\left({\displaystyle \genfrac{}{}{0pt}{}{n+\alpha }{n}}\right)=(1)^n\left({\displaystyle \genfrac{}{}{0pt}{}{\alpha 1}{n}}\right).`$ By this definition, the equality
$$\frac{1}{n!}\left(\frac{d}{dx}\right)^nx^{n+\alpha }|_{x=1}=\left(\genfrac{}{}{0pt}{}{n+\alpha }{n}\right)$$
(4.11)
holds for $`n_0,\alpha `$. Then Laguerre polynomials can be expressed as
$$L_j^\alpha (x)=\underset{\mathrm{}=0}{\overset{j}{}}\frac{(1)^{\mathrm{}}}{\mathrm{}!}\left(\genfrac{}{}{0pt}{}{j+\alpha }{j\mathrm{}}\right)x^{\mathrm{}}$$
(4.12)
for any $`\alpha `$. Remark that applying (4.11) to the equation
$$\frac{1}{n!}\left(\frac{d}{dx}\right)^nx^{n+\alpha }=\frac{1}{(n1)!}\left(\frac{d}{dx}\right)^{n1}x^{(n1)+\alpha }+x\frac{1}{n!}\left(\frac{d}{dx}\right)^nx^{n+(\alpha 1)},$$
with $`x=1`$ and putting $`\beta =\alpha +n`$, we have the identity
$$\left(\genfrac{}{}{0pt}{}{\beta }{n}\right)=\left(\genfrac{}{}{0pt}{}{\beta 1}{n}\right)+\left(\genfrac{}{}{0pt}{}{\beta 1}{n1}\right),n,\beta .$$
(4.13)
We introduce the polynomials
$`F_j(x)`$ $`=`$ $`{\displaystyle \frac{d}{dx}}L_{j+1}^{2𝔞}(x),j_0,`$ (4.14)
$`G_j(x)`$ $`=`$ $`{\displaystyle \frac{d}{dx}}\left\{L_{j+1}^{2𝔞}(x){\displaystyle \frac{j+2𝔞}{j}}L_{j1}^{2𝔞}(x)\right\},j.`$ (4.15)
For $`k_0,j=0,1,2,\mathrm{},k`$, let
$`\alpha _{k,j}`$ $`=`$ $`\left({\displaystyle \genfrac{}{}{0pt}{}{kj+𝔟}{kj}}\right),\text{ if }k\text{ is even},`$
$`\alpha _{k,j}`$ $`=`$ $`{\displaystyle \frac{k+2𝔞}{k}}\left({\displaystyle \genfrac{}{}{0pt}{}{k2j+𝔟}{k2j}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{kj+𝔟}{kj}}\right),\text{ if }k\text{ is odd},`$ (4.16)
In Appendix B, we will give the proof of the following lemmas.
###### Lemma 4.1
For $`\mathrm{}_0`$
$`F_2\mathrm{}(x)={\displaystyle \underset{j=0}{\overset{2\mathrm{}}{}}}\alpha _{2\mathrm{},j}L_j^\nu (x),`$ (4.17)
$`G_{2\mathrm{}+1}(x)={\displaystyle \underset{j=0}{\overset{2\mathrm{}+1}{}}}\alpha _{2\mathrm{}+1,j}L_j^\nu (x).`$ (4.18)
###### Lemma 4.2
For $`q,\mathrm{}_0`$
$`F_{2q},G_{2\mathrm{}+1}_{}=G_{2\mathrm{}+1},F_{2q}_{}=r_q^{}\delta _q\mathrm{},`$ (4.19)
$`F_{2q},F_2\mathrm{}_{}=0,`$ (4.20)
$`G_{2q+1},G_{2\mathrm{}+1}_{}=0,`$ (4.21)
with
$$r_q^{}\frac{4\mathrm{\Gamma }(2q+2𝔞+2)}{(2q+1)!}=4\mathrm{\Gamma }(2𝔞+1)\left(\genfrac{}{}{0pt}{}{2q+2𝔞+1}{2q+1}\right).$$
(4.22)
Then if we define the monic polynomials in $`x`$ of degree $`k`$ for $`k_0`$ as
$$R_k(x)=k!\left(\frac{c_1}{\chi _1}\right)^k\underset{j=0}{\overset{k}{}}\alpha _{k,j}L_j^\nu \left(\frac{x}{c_1}\right)\chi _1^j,$$
(4.23)
Lemma 4.2 gives the following through the relation (4.9) and the orthogonality of the Laguerre polynomials (B.1).
###### Lemma 4.3
For $`q,\mathrm{}_0`$
$`R_{2q},R_{2\mathrm{}+1}=R_{2\mathrm{}+1},R_{2q}=r_q\delta _q\mathrm{},`$
$`R_{2q},R_2\mathrm{}=0,R_{2q+1},R_{2\mathrm{}+1}=0,`$
where
$$r_q=2^{2\nu }T^\kappa \left(\frac{t_1^2}{T}\right)^{4q+1}\frac{(2q)!\mathrm{\Gamma }(2q+2+2𝔞)}{\mathrm{\Gamma }(\nu +1)^2}.$$
(4.24)
The choice of the polynomials $`F_j(x)`$ and $`G_j(x)`$ in (4.14) and (4.15), and their explicit expansions in terms of the Laguerre polynomials (Lemma 4.1) are crucial, since they enable us to determine the appropriate skew-orthogonal polynomials (Lemma 4.3). As shown below, we are able to inverse the skew-symmetric matrix $`A_0`$ given by (4.1) readily for arbitrary (even) $`N`$, by using these skew-orthogonal polynomials.
### 4.3 Matrix inversion
Let $`b_{2k}=b_{2k+1}=r_k^{1/2},k_0`$, and determine the polynomials $`\{M_i(x)\}_{0iN1}`$ in (3.12) as
$$M_i(x)=b_iR_i(x),i=0,1,\mathrm{},N1.$$
Then by (4.1), (4.4) and Lemma 4.3, we have the equality
$$i|\widehat{J}_N|j=(J_N)_{ij},i,j=1,2,\mathrm{},N,$$
(4.25)
where $`J_N=I_{N/2}J_2`$. It is interesting to compare this result with (3.7). Since $`J_N^2=I_N`$, we can immediately obtain the inversion matrix appearing in (3.15) as
$$i|(\widehat{J}_N)^{}|j=(J_N)_{ij},i,j=1,2,\mathrm{},N.$$
(4.26)
If we consider a semi-infinite matrix
$$J\underset{N\mathrm{}}{lim}J_N=\left(i|\widehat{J}|j\right)_{i,j},$$
its inverse matrix may be given by
$$J^1=\left(i|\widehat{J}^{}|j\right)_{i,j}=J.$$
Using expansions (5.6) and (5.8) with Lemmas 4.1, 4.2 and C.1, we can show
$$m,x|\widehat{p}|n,y=m,x|\widehat{p}\widehat{J}|ii|\widehat{J}^{}|jj|n,y$$
and so
$$m,x|\widehat{p}\widehat{J}\widehat{p}|n,y=m,x|\widehat{p}\widehat{J}|ii|\widehat{J}^{}|jj|\widehat{J}\widehat{p}|n,y.$$
Then the equations (3.16) are written as
$`\stackrel{~}{S}^{m,n}(x,y)`$ $`=`$ $`\{\begin{array}{cc}m,x|\widehat{p}\widehat{J}|ii|(\widehat{J}_N)^{}|jj|n,y,\hfill & \text{if }mn\text{,}\hfill \\ & \\ m,x|\widehat{p}\widehat{J}|ii|(\widehat{J}^{}(\widehat{J}_N)^{})|jj|n,y,\hfill & \text{if }m<n\text{,}\hfill \end{array}`$ (4.30)
$`\stackrel{~}{I}^{m,n}(x,y)`$ $`=`$ $`m,x|\widehat{p}\widehat{J}|ii|(\widehat{J}^{}(\widehat{J}_N)^{})|jj|\widehat{J}\widehat{p}|n,y.`$ (4.31)
Now we introduce the notations, just following the previous papers for multi-matrix models , as
$`R_i^{(m)}(x)`$ $``$ $`{\displaystyle \frac{1}{b_i}}m,x|i+1`$ (4.32)
$`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑yR_i(y)\stackrel{~}{p}^{(\nu ,\kappa )}(t_1,y|0)\stackrel{~}{p}^{(\nu ,\kappa )}(t_mt_1,x|y),`$
$`\mathrm{\Phi }_i^{(m)}(x)`$ $``$ $`{\displaystyle \frac{1}{b_i}}m,x|\widehat{p}\widehat{J}|i+1`$ (4.33)
$`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑yR_i^{(m)}(y)F^{(m)}(y,x),`$
for $`i=0,1,\mathrm{},N1,m=1,2,\mathrm{},M+1`$, where
$$F^{(m)}(x,y)=_0^{\mathrm{}}𝑑w_0^w𝑑z\left|\begin{array}{cc}\stackrel{~}{p}^{(\nu ,\kappa )}(Tt_m,z|x)& \stackrel{~}{p}^{(\nu ,\kappa )}(Tt_m,w|x)\\ \stackrel{~}{p}^{(\nu ,\kappa )}(Tt_m,z|y)& \stackrel{~}{p}^{(\nu ,\kappa )}(Tt_m,w|y)\end{array}\right|.$$
(4.34)
It should be noted that $`R_i^{(1)}(x)=R_i(x)\stackrel{~}{p}^{(\nu ,\kappa )}(t_1,x|0),0iN1`$, and $`F^{(1)}(x,y)=F(x,y)`$, where $`R_i(x)`$ and $`F(x,y)`$ were defined by (4.23) and (4.3), respectively. Then we arrive at the following explicit expressions for the elements of matrix kernel (3.23) of our Pfaffian processes,
$`D^{m,n}(x,y)=D_N^{m,n}(x,y)={\displaystyle \underset{\mathrm{}=0}{\overset{(N/2)1}{}}}{\displaystyle \frac{1}{r_{\mathrm{}}}}\left[R_2\mathrm{}^{(m)}(x)R_{2\mathrm{}+1}^{(n)}(y)R_{2\mathrm{}+1}^{(m)}(x)R_2\mathrm{}^{(n)}(y)\right],`$
$`\stackrel{~}{I}^{m,n}(x,y)=\stackrel{~}{I}_N^{m,n}(x,y)={\displaystyle \underset{\mathrm{}=N/2}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{r_{\mathrm{}}}}\left[\mathrm{\Phi }_2\mathrm{}^{(m)}(x)\mathrm{\Phi }_{2\mathrm{}+1}^{(n)}(y)\mathrm{\Phi }_{2\mathrm{}+1}^{(m)}(x)\mathrm{\Phi }_2\mathrm{}^{(n)}(y)\right],`$
$`S^{m,n}(x,y)=S_N^{m,n}(x,y)={\displaystyle \underset{\mathrm{}=0}{\overset{(N/2)1}{}}}{\displaystyle \frac{1}{r_{\mathrm{}}}}\left[\mathrm{\Phi }_2\mathrm{}^{(m)}(x)R_{2\mathrm{}+1}^{(n)}(y)\mathrm{\Phi }_{2\mathrm{}+1}^{(m)}(x)R_2\mathrm{}^{(n)}(y)\right],`$ (4.35)
and
$`\stackrel{~}{S}^{m,n}(x,y)=\stackrel{~}{S}_N^{m,n}(x,y)=S^{m,n}(x,y)\stackrel{~}{p}^{(\nu ,\kappa )}(t_nt_m,y|x)\mathrm{𝟏}_{(m<n)}.`$ (4.36)
## 5 Asymptotic Behavior of Correlation Functions
In this section, we give the proof of our main theorem (Theorem 2.1), by estimating the $`N\mathrm{}`$ asymptotic of matrix kernel (3.23) of Theorem 3.1. Elementary calculation needed for the estimation are summarized in Appendix D. Here $`a_Nb_N,N\mathrm{}`$ means $`a_N/b_N1,N\mathrm{}`$. We assume that $`T=N`$, $`t_m=T+s_m,1mM+1`$ with $`s_1<s_2<\mathrm{}<s_M<s_{M+1}=0`$. We put
$`L_j^\nu (x,s_m)=L_j^\nu \left(x\right)\chi _m^j.\text{ and }\widehat{L}_j^\nu (x,s_m)={\displaystyle \frac{\mathrm{\Gamma }(j+1)}{\mathrm{\Gamma }(j+1+\nu )}}L_j^\nu \left(x\right)\chi _m^j.`$ (5.1)
### 5.1 Asymptotics of $`R_k(x)`$ and $`R_k^{(m)}(x)`$
Let
$$\widehat{R}_k(x)=\frac{1}{k!}\left(\frac{t_1^2}{T}\right)^kR_k(x)=\underset{j=0}{\overset{k}{}}\alpha _{k,j}L_j^\nu (\frac{x}{c_1},s_1).$$
Since $`c_1N=T`$,
$`\widehat{R}_2\mathrm{}(x)`$ $``$ $`I(2\mathrm{},𝔟),`$
$`\widehat{R}_{2\mathrm{}+1}(x)`$ $``$ $`{\displaystyle \frac{2𝔞}{2\mathrm{}+1}}I(2\mathrm{}1,𝔟)I(2\mathrm{}+1,𝔟1)I(2\mathrm{},𝔟1)`$ (5.2)
$``$ $`{\displaystyle \frac{𝔞}{\mathrm{}}}I(2\mathrm{},𝔟)2I(2\mathrm{},𝔟1),N\mathrm{},`$
where
$$I(q,c)\underset{j=0}{\overset{q}{}}\left(\genfrac{}{}{0pt}{}{qj+c}{qj}\right)L_j^\nu (\frac{x}{N},s_1)=\underset{j=0}{\overset{q}{}}\left(\genfrac{}{}{0pt}{}{j+c}{j}\right)L_{qj}^\nu (\frac{x}{N},s_1)$$
for $`q`$ and $`c`$. We set
$$2\mathrm{}=N\theta ,$$
and examine the asymptotic behavior of $`I(2\mathrm{},c)`$ as $`N\mathrm{}`$ with some $`\theta (0,\mathrm{})`$. When $`c_{}`$, $`\left({\displaystyle \genfrac{}{}{0pt}{}{j+c}{j}}\right)=(1)^j\left({\displaystyle \genfrac{}{}{0pt}{}{c1}{j}}\right)`$. Then from (D.10) in Lemma D.2 with $`j=2\mathrm{}`$ (i.e. $`\eta =1`$ in (D.5)), we can easily see
$`I(2\mathrm{},c)`$ $`=`$ $`{\displaystyle \underset{j=0}{\overset{c1}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{j+c}{j}}\right)L_{2\mathrm{}j}^\nu ({\displaystyle \frac{x}{N}},s_1){\displaystyle \frac{(N\theta )^{c+\nu +1}}{(\theta x)^\nu }}\stackrel{~}{J}_\nu ^{(c1)}(\theta ,1,x,s_1),N\mathrm{}.`$
This result is generalized to the following lemma.
###### Lemma 5.1
For any $`c`$, $`\theta (0,\mathrm{})`$, we have
$`I(2\mathrm{},c)`$ $``$ $`{\displaystyle \frac{(N\theta )^{c+\nu +1}}{(\theta x)^\nu }}\stackrel{~}{J}_\nu ^{(c1)}(\theta ,1,x,s_1),N\mathrm{}.`$ (5.3)
Proof.
$`I(2\mathrm{},c)={\displaystyle \underset{p=0}{\overset{2\mathrm{}}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{p+c}{p}}\right)L_{2\mathrm{}p}^\nu ({\displaystyle \frac{x}{N}},s_1)`$
$`={\displaystyle \underset{p=0}{\overset{2\mathrm{}}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{p+c}{p}}\right)\left\{L_{2\mathrm{}p}^\nu ({\displaystyle \frac{x}{N}},s_1)L_2\mathrm{}^\nu ({\displaystyle \frac{x}{N}},s_1)\right\}+{\displaystyle \underset{p=0}{\overset{2\mathrm{}}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{p+c}{p}}\right)L_2\mathrm{}^\nu ({\displaystyle \frac{x}{N}},s_1)`$
$`={\displaystyle \underset{p=0}{\overset{2\mathrm{}}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{p+c}{p}}\right){\displaystyle \underset{k=0}{\overset{p1}{}}}\left\{L_{2\mathrm{}k1}^\nu ({\displaystyle \frac{x}{N}},s_1)L_{2\mathrm{}k}^\nu ({\displaystyle \frac{x}{N}},s_1)\right\}+{\displaystyle \underset{p=0}{\overset{2\mathrm{}}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{p+c}{p}}\right)L_2\mathrm{}^\nu ({\displaystyle \frac{x}{N}},s_1)`$
$`={\displaystyle \underset{p=0}{\overset{2\mathrm{}}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{p+c}{p}}\right){\displaystyle \underset{k=0}{\overset{p1}{}}}{\displaystyle \underset{q=0}{\overset{1}{}}}(1)^{q+1}\left({\displaystyle \genfrac{}{}{0pt}{}{1}{q}}\right)L_{2\mathrm{}kq}^\nu ({\displaystyle \frac{x}{N}},s_1)+{\displaystyle \underset{p=0}{\overset{2\mathrm{}}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{p+c}{p}}\right)L_2\mathrm{}^\nu ({\displaystyle \frac{x}{N}},s_1).`$
Repeating this procedure, we have
$$I(2\mathrm{},c)=\underset{k=0}{\overset{\mathrm{}}{}}(1)^ka_k(2\mathrm{},c)\underset{q=0}{\overset{k}{}}(1)^q\left(\genfrac{}{}{0pt}{}{k}{q}\right)L_{2\mathrm{}q}^\nu (\frac{x}{N},s_1)$$
(5.4)
with
$`a_k(2\mathrm{},c)={\displaystyle \underset{p=0}{\overset{2\mathrm{}}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{p+c}{p}}\right){\displaystyle \underset{j_1=0}{\overset{p1}{}}}{\displaystyle \underset{j_2=0}{\overset{j_11}{}}}\mathrm{}{\displaystyle \underset{j_k=0}{\overset{j_{k1}1}{}}}1={\displaystyle \underset{p=0}{\overset{2\mathrm{}}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{p+c}{p}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{p}{k}}\right).`$
Using (4.13), we can rewrite $`a_k(2\mathrm{},c)`$ as
$`a_k(2\mathrm{},c)={\displaystyle \underset{p=0}{\overset{2\mathrm{}}{}}}\left\{\left({\displaystyle \genfrac{}{}{0pt}{}{p+c+1}{p}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{p+c}{p1}}\right)\right\}\left({\displaystyle \genfrac{}{}{0pt}{}{p}{k}}\right)`$
$`={\displaystyle \underset{p=0}{\overset{2\mathrm{}}{}}}\left[\left({\displaystyle \genfrac{}{}{0pt}{}{p+c+1}{p}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{p}{k}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{p+c}{p1}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{p1}{k}}\right)+\left({\displaystyle \genfrac{}{}{0pt}{}{p+c}{p1}}\right)\left\{\left({\displaystyle \genfrac{}{}{0pt}{}{p1}{k}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{p}{k}}\right)\right\}\right]`$
$`=\left({\displaystyle \genfrac{}{}{0pt}{}{2\mathrm{}+c+1}{2\mathrm{}}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{2\mathrm{}}{k}}\right)a_{k1}(2\mathrm{}1,c+1).`$
Using this equation recursively, we obtain
$$a_k(2\mathrm{},c)=\underset{r=0}{\overset{k}{}}(1)^r\left(\genfrac{}{}{0pt}{}{2\mathrm{}+c+1}{2\mathrm{}r}\right)\left(\genfrac{}{}{0pt}{}{2\mathrm{}r}{kr}\right)=\left(\genfrac{}{}{0pt}{}{2\mathrm{}+c+1}{2\mathrm{}k}\right)\left(\genfrac{}{}{0pt}{}{c+k}{k}\right).$$
(5.5)
Thus (5.4) with (5.5) gives
$$I(2\mathrm{},c)=\underset{k=0}{\overset{\mathrm{}}{}}(1)^k\left(\genfrac{}{}{0pt}{}{2\mathrm{}+c+1}{2\mathrm{}k}\right)\left(\genfrac{}{}{0pt}{}{c+k}{k}\right)\underset{q=0}{\overset{k}{}}(1)^q\left(\genfrac{}{}{0pt}{}{k}{q}\right)L_{2\mathrm{}q}^\nu (\frac{x}{N},s_1).$$
By simple calculation with the estimate (D.2) and (D.10) of Lemma D.2, we obtain (5.3) through the expression (2.16). ∎
From above asymptotic of $`I(2\mathrm{},c)`$ with equations (5.2) we have the following proposition.
###### Proposition 5.2
(1) Suppose that $`\mathrm{}`$ and $`2\mathrm{}N\theta `$, $`N\mathrm{}`$ for some $`\theta (0,\mathrm{})`$. Then
$$\widehat{R}_2\mathrm{}(x)\frac{(N\theta )^{2𝔞+1}}{(\theta x)^\nu }\stackrel{~}{J}_\nu ^{(𝔟1)}(\theta ,1,x,s_1),N\mathrm{}.$$
(2) Suppose that $`\mathrm{}_0`$ and $`2\mathrm{}+1N\theta `$, $`N\mathrm{}`$ for some $`\theta (0,\mathrm{})`$. Then
$$\widehat{R}_{2\mathrm{}+1}(x)\frac{2(N\theta )^{𝔟+\nu }}{(\theta x)^\nu }\left[𝔞\stackrel{~}{J}_\nu ^{(𝔟1)}(\theta ,1,x,s_1)\stackrel{~}{J}_\nu ^{(𝔟)}(\theta ,1,x,s_1)\right],N\mathrm{}.$$
We next examine asymptotic of $`R_k^{(m)}(x)`$. From the definition (4.32) and the expression (4.7)
$`R_k^{(m)}(x)=c_1{\displaystyle _0^{\mathrm{}}}𝑑\eta R_k(c_1\eta )\stackrel{~}{p}^{(\nu ,\kappa )}(t_1,c_1\eta |0)\stackrel{~}{p}^{(\nu ,\kappa )}(t_mt_1,x|c_1\eta )`$
$`={\displaystyle \frac{k!}{2^{\nu +1}\mathrm{\Gamma }(\nu +1)}}\left({\displaystyle \frac{c_1}{\chi _1}}\right)^k\left({\displaystyle \frac{1}{t_m}}\right)^{\nu +1}x^𝔞\mathrm{exp}\left[\left(1+{\displaystyle \frac{t_m}{2T}}\right){\displaystyle \frac{x}{c_m}}\right]{\displaystyle \underset{j=0}{\overset{k}{}}}\alpha _{k,j}L_j^\nu ({\displaystyle \frac{x}{c_m}},s_m).`$
(5.6)
We put
$$\widehat{R}_k^{(m)}(x)=\frac{2^\nu T^\nu \mathrm{\Gamma }(\nu +1)}{\mathrm{\Gamma }(k+1+2𝔞)}\left(\frac{\chi _1}{c_1}\right)^kR_k^{(m)}(x).$$
(5.7)
If we set $`kN\theta `$ as $`N\mathrm{}`$, (D.1) in Appendix D gives
$$\frac{k!T^\nu }{\mathrm{\Gamma }(k+1+2𝔞)}\left(\frac{1}{t_m}\right)^{\nu +1}\mathrm{exp}\left[\left(1+\frac{t_m}{2T}\right)\frac{x}{c_m}\right]N^{(2𝔞+1)}\theta ^{2𝔞},N\mathrm{},$$
then we obtain the following from Proposition 5.2.
###### Proposition 5.3
(1) Suppose that $`\mathrm{}`$ and $`2\mathrm{}N\theta `$, $`N\mathrm{}`$ for some $`\theta (0,\mathrm{})`$. Then
$$\widehat{R}_2\mathrm{}^{(m)}(x)\frac{\theta ^{1\nu }x^{\kappa /2}}{2}\stackrel{~}{J}_\nu ^{(𝔟1)}(\theta ,1,x,s_m),N\mathrm{}.$$
(2) Suppose that $`\mathrm{}_0`$ and $`2\mathrm{}+1N\theta `$, $`N\mathrm{}`$ for some $`\theta (0,\mathrm{})`$. Then
$$\widehat{R}_{2\mathrm{}+1}^{(m)}(x)\frac{\theta ^\nu x^{\kappa /2}}{N}\left[𝔞\stackrel{~}{J}_\nu ^{(𝔟1)}(\theta ,1,x,s_m)\stackrel{~}{J}_\nu ^{(𝔟)}(\theta ,1,x,s_m)\right],N\mathrm{}.$$
### 5.2 Asymptotics of $`\mathrm{\Phi }_k^{(m)}(x)`$
Using (4.8), (4.34) is rewritten as
$`F^{(m)}(y,x)`$ $`=`$ $`\left({\displaystyle \frac{1}{T}}\right)^\kappa \left({\displaystyle \frac{t_m}{c_m}}\right)^{2(\nu +1)}(xy)^{\kappa /2}\mathrm{exp}\left\{\left({\displaystyle \frac{t_m}{T}}\right){\displaystyle \frac{x+y}{2c_m}}\right\}`$
$`\times {\displaystyle \underset{p=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}L_p^\nu ,L_j^\nu _{}\widehat{L}_p^\nu ({\displaystyle \frac{y}{c_m}},s_m)\widehat{L}_j^\nu ({\displaystyle \frac{x}{c_m}},s_m).`$
Hence from (4.33) and (5.6), we have
$`\mathrm{\Phi }_k^{(m)}(x)`$ $`=`$ $`c_m{\displaystyle _0^{\mathrm{}}}𝑑\eta R_k^{(m)}(c_m\eta )F^{(m)}(c_m\eta ,x)`$ (5.8)
$`=`$ $`{\displaystyle \frac{k!}{2^{\nu +1}\mathrm{\Gamma }(\nu +1)}}\left({\displaystyle \frac{c_1}{\chi _1}}\right)^k\left({\displaystyle \frac{1}{c_m}}\right)^{\nu +1}{\displaystyle \frac{t_m^{\nu +1}}{T^\kappa }}x^{\kappa /2}\mathrm{exp}\left\{{\displaystyle \frac{t_mx}{2Tc_m}}\right\}`$
$`\times {\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{p=0}{\overset{k}{}}}\alpha _{k,p}L_p^\nu ,L_j^\nu _{}\widehat{L}_j^\nu ({\displaystyle \frac{x}{c_m}},s_m),`$
where we have used the orthogonal relation (B.1) of Laguerre polynomials. Put
$$\widehat{\mathrm{\Phi }}_k^{(m)}(x)=\frac{2^\nu T^𝔟\mathrm{\Gamma }(\nu +1)}{k!}\left(\frac{\chi _1}{c_1}\right)^k\mathrm{\Phi }_k^{(m)}(x).$$
(5.9)
Then we have the following proposition.
###### Proposition 5.4
(1) Suppose that $`\mathrm{}`$ and $`2\mathrm{}N\theta `$, $`N\mathrm{}`$ for some $`\theta (0,\mathrm{})`$. Then
$$\widehat{\mathrm{\Phi }}_2\mathrm{}^{(m)}(x)\theta ^\nu x^{\kappa /2}_1^{\mathrm{}}𝑑\xi \xi ^𝔞\widehat{J}_\nu ^{(𝔟+1)}(\theta ,\xi ,x,s_m),\mathrm{}.$$
(5.10)
(2) Suppose that $`\mathrm{}_0`$ and $`2\mathrm{}+1N\theta `$, $`N\mathrm{}`$ for some $`\theta (0,\mathrm{})`$. Then
$$\widehat{\mathrm{\Phi }}_{2\mathrm{}+1}^{(m)}(x)\frac{2\theta ^{1+\nu }x^{\kappa /2}}{N}\widehat{J}_\nu ^{(𝔟+1)}(\theta ,1,x,s_m),N\mathrm{}.$$
(5.11)
Proof.
$`\widehat{\mathrm{\Phi }}_k^{(m)}(x)`$ $``$ $`{\displaystyle \frac{T^\nu x^{\kappa /2}}{2}}{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{p=0}{\overset{k}{}}}\alpha _{k,p}L_p^\nu ,L_j^\nu _{}\widehat{L}_j^\nu ({\displaystyle \frac{x}{N}},s_m),N\mathrm{},`$
$`=`$ $`{\displaystyle \frac{T^\nu x^{\kappa /2}}{2}}{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}Q_k,{\displaystyle \underset{q=0}{\overset{j}{}}}\beta _{j,q}Q_q_{}\widehat{L}_j^\nu ({\displaystyle \frac{x}{N}},s_m).`$
Here we have used the notation (C.1) and introduced $`𝜷=(\beta _{j,k})`$, which is the inverse of the matrix $`𝜶=(\alpha _{k,j})`$ given by Lemma C.1 in Appendix C. By the skew orthogonality of $`\{Q_k\}`$ given by Lemma 4.2, we have
$`\widehat{\mathrm{\Phi }}_2\mathrm{}^{(m)}(x){\displaystyle \frac{T^\nu x^{\kappa /2}r_{\mathrm{}}^{}}{2}}{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}\beta _{j,2\mathrm{}+1}\widehat{L}_j^\nu ({\displaystyle \frac{x}{N}},s_m),N\mathrm{},`$ (5.12)
$`\widehat{\mathrm{\Phi }}_{2\mathrm{}+1}^{(m)}(x){\displaystyle \frac{T^\nu x^{\kappa /2}r_{\mathrm{}}^{}}{2}}{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}\beta _{j,2\mathrm{}}\widehat{L}_j^\nu ({\displaystyle \frac{x}{N}},s_m),N\mathrm{}.`$ (5.13)
By (4.22) and (C.4), (5.13) gives
$`\widehat{\mathrm{\Phi }}_{2\mathrm{}+1}^{(m)}(x)2\mathrm{\Gamma }(2𝔞+1)T^\nu x^{\kappa /2}\left({\displaystyle \genfrac{}{}{0pt}{}{2\mathrm{}+1+2𝔞}{2\mathrm{}+1}}\right)`$
$`\times {\displaystyle \underset{j=2\mathrm{}}{\overset{\mathrm{}}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{j2\mathrm{}𝔟2}{j2\mathrm{}}}\right)\widehat{L}_j^\nu ({\displaystyle \frac{x}{N}},s_m),N\mathrm{}.`$ (5.14)
From (D.8) we have
$`{\displaystyle \underset{j=2\mathrm{}}{\overset{\mathrm{}}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{j2\mathrm{}𝔟2}{j2\mathrm{}}}\right)\widehat{L}_j^\nu ({\displaystyle \frac{x}{N}},s_m)={\displaystyle \underset{r=0}{\overset{\mathrm{}}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{r𝔟2}{r}}\right)\widehat{L}_{2\mathrm{}+r}^\nu ({\displaystyle \frac{x}{N}},s_m)`$
$`={\displaystyle \underset{r=0}{\overset{\mathrm{}}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{r𝔟2+\alpha }{r}}\right){\displaystyle \underset{p=0}{\overset{\alpha }{}}}(1)^p\left({\displaystyle \genfrac{}{}{0pt}{}{\alpha }{p}}\right)\widehat{L}_{2\mathrm{}+r+p}^\nu ({\displaystyle \frac{x}{N}},s_m)`$
$`(2\mathrm{})^\alpha {\displaystyle \underset{r=0}{\overset{\mathrm{}}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{r𝔟2+\alpha }{r}}\right)\widehat{J}_\nu ^{(\alpha )}(\theta ,\eta ,x,s_m),N\mathrm{},`$ (5.15)
where (D.11) of Lemma D.2 was applied. Setting $`j=2\mathrm{}\eta =N\theta \eta `$ and using (D.2) in Appendix D, we conclude that
$`\widehat{\mathrm{\Phi }}_{2\mathrm{}+1}^{(m)}(x)`$ $``$ $`2T^\nu (2\mathrm{})^{2𝔞\alpha }x^{\kappa /2}{\displaystyle \underset{j=2\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(j2\mathrm{}+1)^{𝔟2+\alpha }}{\mathrm{\Gamma }(𝔟1+\alpha )}}\widehat{J}_\nu ^{(\alpha )}(\theta ,\eta ,x,s_m)`$ (5.16)
$``$ $`{\displaystyle \frac{2\theta ^{1+\nu }x^{\kappa /2}}{N\mathrm{\Gamma }(𝔟1+\alpha )}}{\displaystyle _1^{\mathrm{}}}𝑑\eta {\displaystyle \frac{\widehat{J}_\nu ^{(\alpha )}(\theta ,\eta ,x,s_m)}{(\eta 1)^{𝔟+2\alpha }}},N\mathrm{}.`$
Through the expression (2.17), we obtain (5.11).
By (C.5) of Lemma C.1 with (C.2)
$`{\displaystyle \underset{j=2\mathrm{}+1}{\overset{\mathrm{}}{}}}\beta _{j,2\mathrm{}+1}\widehat{L}_j^\nu ({\displaystyle \frac{x}{N}},s_m)`$
$`={\displaystyle \frac{1}{b(1,2\mathrm{}+1)}}{\displaystyle \underset{j=2\mathrm{}+1}{\overset{\mathrm{}}{}}}\widehat{L}_j^\nu ({\displaystyle \frac{x}{N}},s_m){\displaystyle \underset{r=\mathrm{}+1}{\overset{[(j+1)/2]}{}}}b(1,2r1)\left({\displaystyle \genfrac{}{}{0pt}{}{j2r𝔟1}{j2r+1}}\right)`$
$`={\displaystyle \frac{1}{b(1,2\mathrm{}+1)}}S(\mathrm{}),`$
where
$$S(\mathrm{})=\underset{r=\mathrm{}+1}{\overset{\mathrm{}}{}}b(1,2r1)\underset{j=2r1}{\overset{\mathrm{}}{}}\widehat{L}_j^\nu (\frac{x}{N},s_m)\left(\genfrac{}{}{0pt}{}{j2r𝔟1}{j2r+1}\right).$$
By this equation with the estimate (D.3) for $`b(1,2r1)`$ and (4.22), (5.12) becomes
$`\widehat{\mathrm{\Phi }}_2\mathrm{}^{(m)}(x)`$ $``$ $`2T^\nu (2\mathrm{}+2)^𝔞\mathrm{\Gamma }(2𝔞+1)\left({\displaystyle \genfrac{}{}{0pt}{}{2\mathrm{}+1+2𝔞}{2\mathrm{}+1}}\right)x^{\kappa /2}S(\mathrm{})`$ (5.17)
$``$ $`2T^\nu (2\mathrm{}+2)^𝔞x^{\kappa /2}S(\mathrm{}),N\mathrm{}.`$
From (5.15) with (D.3)
$`S(\mathrm{})`$ $``$ $`{\displaystyle \underset{r=\mathrm{}+1}{\overset{\mathrm{}}{}}}(2r)^𝔞\left({\displaystyle \frac{1}{2\mathrm{}}}\right)^\alpha {\displaystyle \underset{j=2r1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(j2r+1)^{𝔟2+\alpha }}{\mathrm{\Gamma }(𝔟1+\alpha )}}\widehat{J}_\nu ^{(\alpha )}(\theta ,\eta ,x,s_m)`$
$``$ $`{\displaystyle \frac{(2\mathrm{})^{\kappa /2}}{2\mathrm{\Gamma }(𝔟1+\alpha )}}{\displaystyle _1^{\mathrm{}}}𝑑\xi \xi ^𝔞{\displaystyle _\xi ^{\mathrm{}}}𝑑\eta {\displaystyle \frac{\widehat{J}_\nu ^{(\alpha )}(\theta ,\eta ,x,s_m)}{(\eta \xi )^{𝔟+2\alpha }}},N\mathrm{}.`$
Thus
$`\widehat{\mathrm{\Phi }}_2\mathrm{}^{(m)}(x)`$ $``$ $`{\displaystyle \frac{\theta ^\nu x^{\kappa /2}}{\mathrm{\Gamma }(𝔟1+\alpha )}}{\displaystyle _1^{\mathrm{}}}𝑑\xi \xi ^𝔞{\displaystyle _\xi ^{\mathrm{}}}𝑑\eta {\displaystyle \frac{\widehat{J}_\nu ^{(\alpha )}(\theta ,\eta ,x,s_m)}{(\eta \xi )^{𝔟+2\alpha }}}`$
$`=`$ $`\theta ^\nu x^{\kappa /2}{\displaystyle _1^{\mathrm{}}}𝑑\xi \xi ^𝔞\widehat{J}_\nu ^{(𝔟+1)}(\theta ,\xi ,x,s_m),N\mathrm{},`$
where we used the expression (2.17). This completes the proof of Proposition 5.4. ∎
### 5.3 Asymptotics of $`D^{m,n}(x,y)`$, $`\stackrel{~}{I}^{m,n}(x,y)`$, $`S^{m,n}(x,y)`$ and $`\stackrel{~}{S}^{m,n}(x,y)`$
From the expressions (4.35) with (4.24) and the definitions (5.7) and (5.9) we have
$`D_N^{m,n}(x,y){\displaystyle \underset{\mathrm{}=0}{\overset{(N/2)1}{}}}\left({\displaystyle \frac{2\mathrm{}}{N}}\right)^{2𝔞}\left[\widehat{R}_2\mathrm{}^{(m)}(x)\widehat{R}_{2\mathrm{}+1}^{(n)}(y)\widehat{R}_{2\mathrm{}+1}^{(m)}(x)\widehat{R}_2\mathrm{}^{(n)}(y)\right],`$
$`\stackrel{~}{I}_N^{m,n}(x,y){\displaystyle \underset{\mathrm{}=(N/2)}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{2\mathrm{}}{N}}\right)^{2𝔞}\left[\widehat{\mathrm{\Phi }}_2\mathrm{}^{(m)}(x)\widehat{\mathrm{\Phi }}_{2\mathrm{}+1}^{(n)}(y)\widehat{\mathrm{\Phi }}_{2\mathrm{}+1}^{(m)}(x)\widehat{\mathrm{\Phi }}_2\mathrm{}^{(n)}(y)\right],`$
$`S_N^{m,n}(x,y){\displaystyle \underset{\mathrm{}=0}{\overset{(N/2)1}{}}}\left[\widehat{\mathrm{\Phi }}_2\mathrm{}^{(m)}(x)\widehat{R}_{2\mathrm{}+1}^{(n)}(y)\widehat{\mathrm{\Phi }}_{2\mathrm{}+1}^{(m)}(x)\widehat{R}_2\mathrm{}^{(n)}(y)\right],N\mathrm{}.`$
From Propositions 5.3 and 5.4 we obtain the following asymptotics:
$`D_N^{m,n}(x,y)`$ $``$ $`𝒟(s_m,x;s_n,y),`$
$`\stackrel{~}{I}_N^{m,n}(x,y)`$ $``$ $`\stackrel{~}{}(s_m,x;s_n,y),`$
$`S_N^{m,n}(x,y)`$ $``$ $`𝒮(s_m,x;s_n,y),N\mathrm{},`$
where $`𝒟,\stackrel{~}{},𝒮`$ are defined by (2.19).
Next we study the asymptotic behavior of $`\stackrel{~}{p}^{(\nu ,\kappa )}(t_nt_m,y|x)`$. From (4.7) we have
$`\stackrel{~}{p}^{(\nu ,\kappa )}(t_nt_m,y|x)`$
$`=\left({\displaystyle \frac{t_m}{t_n}}\right)^{\nu +1}c_m^{𝔞1}\left({\displaystyle \frac{x}{c_m}}\right)^{\kappa /2}y^𝔞\mathrm{exp}\left[{\displaystyle \frac{t_mx}{2Tc_m}}\right]\mathrm{exp}\left[\left(2+{\displaystyle \frac{t_n}{T}}\right){\displaystyle \frac{y}{2c_n}}\right]`$
$`\times {\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}\widehat{L}_j^\nu ({\displaystyle \frac{x}{c_m}},s_m)L_j^\nu ({\displaystyle \frac{y}{c_n}},s_n).`$
Then by simple calculation with Lemma D.2 with $`\alpha =0`$, we have
$`\stackrel{~}{p}^{(\nu ,\kappa )}(t_nt_m,y|x)`$ $``$ $`\left({\displaystyle \frac{y}{x}}\right)^{𝔟/2}{\displaystyle \frac{1}{N}}{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}\mathrm{exp}\left[2(s_ms_n)\theta \eta \right]J_\nu (2\sqrt{\theta \eta x})J_\nu (2\sqrt{\theta \eta y}),`$
$``$ $`\left({\displaystyle \frac{y}{x}}\right)^{𝔟/2}𝒢(s_m,x;s_n,y),N\mathrm{},`$
where $`𝒢`$ is defined by (2.21), and then $`\stackrel{~}{S}_N^{m,n}(x,y)\stackrel{~}{𝒮}(s_m,x;s_n,y),N\mathrm{}.`$ Then, the proof of Theorem 2.1 is completed.
Appendices
## Appendix A Proof of (3.14)
We assume that the number of particles $`N`$ is even. Consider the multiple integral
$`Z_{N,T}^𝐘[\chi ]`$ $`=`$ $`\left({\displaystyle \frac{1}{N!}}\right)^{M+1}{\displaystyle _{_+^{N(M+1)}}}{\displaystyle \underset{m=1}{\overset{M+1}{}}}d𝐱^{(m)}\underset{1i,jN}{det}\left[M_{i1}(x_j^{(1)})\stackrel{~}{p}^{(\nu ,\kappa )}(t_1,x_j^{(1)}|0)(1+\chi _1(x_j^{(1)}))\right]`$
$`\times {\displaystyle \underset{m=1}{\overset{M}{}}}\underset{1i,jN}{det}\left[\stackrel{~}{p}^{(\nu ,\kappa )}(t_{m+1}t_m,x_j^{(m+1)}|x_i^{(m)})(1+\chi _{m+1}(x_j^{(m+1)}))\right]\mathrm{sgn}\left(h_N(𝐱^{(M+1)})\right).`$
By the definition (3.4) with (3.3) and (3.1), and by the equality (3.11), we have
$$\mathrm{\Psi }_{N,T}^𝐘(𝐟;𝜽)=\frac{Z_{N,T}^𝐘[\chi ]}{Z_{N,T}^𝐘[0]},$$
(A.1)
where $`Z_{N,T}^𝐘[0]`$ is obtained from $`Z_{N,T}^𝐘[\chi ]`$ by setting $`\chi _m(x)0`$ for all $`m=1,2,\mathrm{},M+1`$.
By repeated applications of the Heine identity
$$_{_{+<}^N}𝑑𝐱\underset{1i,jN}{det}\left[\varphi _i(x_j)\right]\underset{1i,jN}{det}\left[\overline{\varphi }_i(x_j)\right]=\underset{1i,jN}{det}\left[__+𝑑x\varphi _i(x)\overline{\varphi }_j(x)\right]$$
for square integrable continuous functions $`\varphi _i,\overline{\varphi }_i,1iN`$, we have
$`Z_{N,T}^𝐘[\chi ]`$ $`=`$ $`{\displaystyle _{_{+<}^N}}d𝐲\underset{1i,jN}{det}[{\displaystyle _{_+^{M+1}}}{\displaystyle \underset{m=1}{\overset{M+1}{}}}dx^{(m)}\left\{M_{i1}(x^{(1)})\stackrel{~}{p}^{(\nu ,\kappa )}(t_1,x^{(1)}|0)(1+\chi _1(x^{(1)}))\right\}`$
$`\times {\displaystyle \underset{m=1}{\overset{M}{}}}\left\{\stackrel{~}{p}^{(\nu ,\kappa )}(t_{m+1}t_m,x^{(m+1)}|x^{(m)})(1+\chi _{m+1}(x^{(m+1)}))\right\}\delta (y_jx^{(M+1)})].`$
Using the notations in Section 3.2, it is expressed as
$`Z_{N,T}^𝐘[\chi ]`$ $`=`$ $`{\displaystyle _{_{+<}^N}}𝑑𝐲\underset{1i,jN}{det}\left[i|\left(1+{\displaystyle \frac{1}{1\widehat{\chi }\widehat{p}_+}}\widehat{\chi }\widehat{p}\right)_N|M+1,y_j\right]`$
$`=`$ $`{\displaystyle _{_{+<}^N}}𝑑𝐲\underset{1i,jN}{det}\left[M+1,y_i|\left(1+\widehat{p}\widehat{\chi }{\displaystyle \frac{1}{1\widehat{p}_{}\widehat{\chi }}}\right)_N|j\right],`$
since $`m,x|(\widehat{p}_+)^k|n,y=n,y|(\widehat{p}_{})^k|m,x0`$ for $`k>nm0`$. Here we have used the Chapman-Kolmogorov equation, $`__+𝑑y\stackrel{~}{p}^{(\nu ,\kappa )}(ts,y|x)\stackrel{~}{p}^{(\nu ,\kappa )}(ut,z|y)=\stackrel{~}{p}^{(\nu ,\kappa )}(us,z|x)`$, $`0stuT,x,y_+`$. Next we use the formula of de Bruijin
$$_{_{+<}^N}𝑑𝐲\underset{1i,jN}{det}\left[\varphi _i(y_j)\right]=\mathrm{Pf}_{1i,jN}\left[__+𝑑y__+𝑑\stackrel{~}{y}\mathrm{sgn}(\stackrel{~}{y}y)\varphi _i(y)\varphi _j(\stackrel{~}{y})\right]$$
for integrable continuous functions $`\varphi _i,1iN`$, in which the Pfaffian is defined by (2.22). Since $`(\mathrm{Pf}(A))^2=detA`$ for any even-dimensional skew-symmetric matrix $`A`$, we have
$`\left(Z_{N,T}^𝐘[\chi ]\right)^2`$ $`=`$ $`\underset{1i,jN}{det}\left[i|\left(1+{\displaystyle \frac{1}{1\widehat{\chi }p_+}}\widehat{\chi }\widehat{p}\right)_N\widehat{J}\left(1+\widehat{p}\widehat{\chi }{\displaystyle \frac{1}{1\widehat{p}_{}\widehat{\chi }}}\right)_N|j\right]`$
$`=`$ $`\underset{1i,jN}{det}\left[(A_0)_{ij}+(A_1)_{ij}+(A_2)_{ij}+(A_3)_{ij}\right]`$
with
$`(A_0)_{ij}`$ $`=`$ $`i|\widehat{J}_N|j,`$
$`(A_1)_{ij}`$ $`=`$ $`i|\left({\displaystyle \frac{1}{1\widehat{\chi }\widehat{p}_+}}\widehat{\chi }\widehat{p}\widehat{J}\right)_N|j=i|\left(\widehat{\chi }{\displaystyle \frac{1}{1\widehat{p}_+\widehat{\chi }}}\widehat{p}\widehat{J}\right)_N|j,`$
$`(A_2)_{ij}`$ $`=`$ $`i|\left(\widehat{J}\widehat{p}\widehat{\chi }{\displaystyle \frac{1}{1\widehat{p}_{}\widehat{\chi }}}\right)_N|j,`$
$`(A_3)_{ij}`$ $`=`$ $`i|\left(\widehat{\chi }{\displaystyle \frac{1}{1\widehat{p}_+\widehat{\chi }}}\widehat{p}\widehat{J}\widehat{p}\widehat{\chi }{\displaystyle \frac{1}{1\widehat{p}_{}\widehat{\chi }}}\right)_N|j.`$
Since $`\left(Z_{N,T}[0]\right)^2=det_{1i,jN}\left[(A_0)_{ij}\right]`$, (A.1) gives
$$\left\{\mathrm{\Psi }_{N,T}(𝐟;𝜽)\right\}^2=\underset{1i,jN}{det}\left[\delta _{ij}+(A_0^1A_1)_{ij}+(A_0^1A_2)_{ij}+(A_0^1A_3)_{ij}\right].$$
(A.2)
By our notation (3.13), $`(A_0^1)_{ij}=i|(\widehat{J}_N)^{}|j,`$ and it is easy to confirm that (A.2) is written in the form
$`\left\{\mathrm{\Psi }_{N,T}^𝐘(𝐟;𝜽)\right\}^2`$ $`=`$ $`\underset{1i,jN}{det}\left[\delta _{ij}+i|𝐁|m,xm,x|𝐂|j\right],`$ (A.3)
where we have introduced $`𝐁`$ and $`𝐂`$ as the following two-dimensional row and column vector-valued operators,
$`𝐁`$ $`=`$ $`\left(\begin{array}{ccc}(\widehat{J}_N)^{}\widehat{\chi }\frac{1}{1\widehat{p}_+\widehat{\chi }}\hfill & & (\widehat{J}_N)^{}\left(1+\widehat{\chi }\frac{1}{1\widehat{p}_+\widehat{\chi }}\widehat{p}\right)\widehat{J}\widehat{p}\widehat{\chi }\hfill \end{array}\right),`$ (A.5)
$`𝐂`$ $`=`$ $`\left(\begin{array}{c}\widehat{p}\widehat{J}\\ \\ \frac{1}{1\widehat{p}_{}\widehat{\chi }}\end{array}\right).`$ (A.9)
The determinant (A.3) is equivalent with the Fredholm determinant,
$$\mathrm{Det}m,x|I_2+𝐂𝐁|n,y.$$
Introducing matrix-valued operators,
$$𝐊_+=\left(\begin{array}{ccc}1\widehat{p}_+\widehat{\chi }& & 0\\ & & \\ 0& & 1\end{array}\right),𝐊_{}=\left(\begin{array}{ccc}1& & 0\\ & & \\ 0& & 1\widehat{p}_{}\widehat{\chi }\end{array}\right),\widehat{𝐊}=\left(\begin{array}{ccc}1& & \widehat{p}\widehat{J}\widehat{p}\widehat{\chi }\\ & & \\ 0& & 1\end{array}\right),$$
we have
$`I_2+𝐂𝐁`$ $`=`$ $`𝐊_{}^{}{}_{}{}^{1}\left[𝐊_{}𝐊_++\left(\begin{array}{ccc}\widehat{p}\widehat{J}𝒥_N^{}& & \widehat{p}\widehat{J}\left(1𝒥_N^{}\widehat{J}\right)\widehat{p}\\ & & \\ 𝒥_N^{}& & 𝒥_N^{}\widehat{J}\widehat{p}\end{array}\right)\widehat{\chi }\right]𝐊_{+}^{}{}_{}{}^{1}\widehat{𝐊}`$ (A.13)
$`=`$ $`𝐊_{}^{}{}_{}{}^{1}\left[I_2+\left(\begin{array}{ccc}\widehat{p}\widehat{J}𝒥_N^{}\widehat{p}_+& & \widehat{p}\widehat{J}\widehat{p}\widehat{p}\widehat{J}𝒥_N^{}\widehat{J}\widehat{p}\\ & & \\ 𝒥_N^{}& & 𝒥_N^{}\widehat{J}\widehat{p}\widehat{p}_{}\end{array}\right)\widehat{\chi }\right]𝐊_{+}^{}{}_{}{}^{1}\widehat{𝐊},`$ (A.17)
where $`𝒥_N^{}=(\widehat{J}_N)^{}`$. From the orthogonality (3.6) and the definitions (3.9) of the operators $`\widehat{p}_+`$ and $`\widehat{p}_{}`$, we have the fact that
$$\mathrm{Det}m,x|𝐊_+|n,y=\mathrm{Det}m,x|𝐊_{}|n,y=\mathrm{Det}m,x|\widehat{𝐊}|n,y=1.$$
Then (3.14) is derived.
## Appendix B Proofs of Lemmas 4.1 and 4.2
We use the following orthogonal relations and formulae on Laguerre polynomials, which hold for $`\alpha ,\beta >1`$;
$`{\displaystyle _0^{\mathrm{}}}L_j^\alpha (x)L_k^\alpha (x)x^\alpha e^x𝑑x={\displaystyle \frac{\mathrm{\Gamma }(\alpha +j+1)}{\mathrm{\Gamma }(j+1)}}\delta _{jk},j,k_0,`$ (B.1)
$`x{\displaystyle \frac{d}{dx}}L_j^\alpha (x)=jL_j^\alpha (x)(j+\alpha )L_{j1}^\alpha (x),j,`$ (B.2)
$`L_j^\alpha (x)={\displaystyle \frac{d}{dx}}L_{j+1}^\alpha (x)+{\displaystyle \frac{d}{dx}}L_j^\alpha (x),j_0,`$ (B.3)
$`L_j^\beta (x)={\displaystyle \underset{k=0}{\overset{j}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{jk+\beta \alpha 1}{jk}}\right)L_k^\alpha (x),j_0.`$ (B.4)
Remark 6. The identities (B.2) and (B.3) are given as Eqs. (6.2.6) and (6.2.7) in . The relation (B.4) is proved in as (6.2.37) only when $`\beta \alpha >1`$. The identity (see (4.10) and )
$$\underset{\mathrm{}=0}{\overset{k}{}}\left(\genfrac{}{}{0pt}{}{\mathrm{}\alpha 1}{\mathrm{}}\right)\left(\genfrac{}{}{0pt}{}{k\mathrm{}+\alpha 1}{k\mathrm{}}\right)=\delta _{k0},$$
can be used to invert the relation (B.4) to the form
$$L_{\mathrm{}}^\alpha (x)=\underset{j=0}{\overset{\mathrm{}}{}}\left(\genfrac{}{}{0pt}{}{\mathrm{}j+\alpha \beta 1}{\mathrm{}j}\right)L_j^\beta (x),\mathrm{}_0.$$
Therefore, the validity of (B.4) for $`\beta \alpha >1`$ implies that for $`\alpha >\beta >1`$.
### B.1 Proof of Lemma 4.1
In this subsection we prove Lemma 4.1, which gives the expansion formulae of $`F_k(x)`$ and $`G_k(x)`$ in terms of $`\{L_j^\nu (x)\}`$. Taking the summation of the equalities (B.3) from $`0`$ to $`k_0`$ gives
$$\underset{n=0}{\overset{k}{}}L_n^{2𝔞}(x)=\frac{d}{dx}L_{k+1}^{2𝔞}(x)+\frac{d}{dx}L_0^{2𝔞}(x)=F_k(x).$$
(B.5)
From (B.4), (B.5) and (4.13), we have
$`F_k(x)`$ $`=`$ $`L_k^{2𝔞}(x)+F_{k1}(x)`$
$`=`$ $`{\displaystyle \underset{j=0}{\overset{k}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{kj+𝔟1}{kj}}\right)L_j^\nu (x)+F_{k1}(x)`$
$`=`$ $`{\displaystyle \underset{j=0}{\overset{k}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{kj+𝔟}{kj}}\right)L_j^\nu (x){\displaystyle \underset{j=0}{\overset{k1}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{kj+𝔟1}{kj1}}\right)L_j^\nu (x)+F_{k1}(x).`$
Since $`L_0^\nu (x)=1`$ and $`F_0(x)=L_0^{2𝔞}(x)=1`$,
$`F_k(x){\displaystyle \underset{j=0}{\overset{k}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{kj+𝔟}{kj}}\right)L_j^\nu (x)=F_{k1}(x){\displaystyle \underset{j=0}{\overset{k1}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{k1j+𝔟}{k1j}}\right)L_j^\nu (x)`$
$`=F_0(x)\left({\displaystyle \genfrac{}{}{0pt}{}{𝔟}{0}}\right)L_0^\nu (x)=0.`$
Then we have (4.17). From (4.14) and (4.15) we have
$$G_k(x)=F_k(x)+\frac{k+2𝔞}{k}F_{k2}(x),k.$$
(B.6)
Then (4.17) gives (4.18). This completes the proof. ∎
### B.2 Proof of Lemma 4.2
We introduce a symmetric inner product
$$(f,g)_0^{\mathrm{}}𝑑xe^xx^{2𝔞}f(x)g(x).$$
It is easy to see that it is related with the elementary skew-symmetric inner product (4.5) as
$$f,g_{}=2(f,g)_0^{\mathrm{}}𝑑xe^{x/2}x^𝔞f(x)_0^{\mathrm{}}𝑑ye^{y/2}y^𝔞g(y),$$
(B.7)
where
$$f(x)\frac{_0^x𝑑ze^{z/2}z^𝔞f(z)}{e^{x/2}x^𝔞}.$$
We consider the polynomials
$`W_j(x)`$ $`=`$ $`L_j^{2𝔞}(x){\displaystyle \frac{j+2𝔞}{j}}L_{j1}^{2𝔞}(x),j.`$ (B.8)
###### Lemma B.1
For $`j=_0,k`$
$`(W_k,F_j)={\displaystyle \frac{\mathrm{\Gamma }(j+2𝔞+2)}{(j+1)!}}\delta _{k1j}.`$ (B.9)
Proof. For $`2𝔞>1`$ the above equation can be derived immediately from the definitions (B.8), (B.5), and the orthogonal relation (B.1). We can extend it to the case $`𝔞>1`$. For $`k=2,3,\mathrm{}`$ it is easy to see that
$`(L_k^{2𝔞},x^j)`$ $`=`$ $`(1)^k\mathrm{\Gamma }(k+2𝔞+1)\delta _{jk},j=1,2,\mathrm{},k.`$ (B.10)
However, if $`2𝔞<1`$, $`(L_k^{2𝔞},1)=\mathrm{}`$. Then we use the following equation:
$$(L_k^{2𝔞},L_j^{2𝔞}L_j^{2𝔞}(0))=\frac{\mathrm{\Gamma }(k+2𝔞+1)}{k!}\delta _{jk},k_0,$$
(B.11)
which is obtained from (B.10) for $`j=1,2,\mathrm{},k`$, since $`L_k^{2𝔞}(x)`$ is a polynomial whose coefficient of the $`k`$-th order is $`(1)^k/k!`$. By simple calculation we have
$$(1,L_j^{2𝔞}L_j^{2𝔞}(0))=0,j_0,$$
and thus
$$(L_k^{2𝔞},L_j^{2𝔞}L_j^{2𝔞}(0))=(L_k^{2𝔞}L_k^{2𝔞}(0),L_j^{2𝔞}L_j^{2𝔞}(0))=0,jk+4.$$
From the definitions of $`W_k`$ and $`F_01`$ it is easily to see that
$`(W_k,F_0)`$ $`=`$ $`\mathrm{\Gamma }(2𝔞+2)\delta _{k\mathrm{1\; 0}}.`$ (B.12)
When $`k,j`$, from (B.12) and (B.11),
$`(W_k,F_j)`$ $`=`$ $`(L_k^{2𝔞}{\displaystyle \frac{k+2𝔞}{k}}L_{k1}^{2𝔞},{\displaystyle \underset{p=0}{\overset{j}{}}}L_p^{2𝔞})`$ (B.13)
$`=`$ $`{\displaystyle \underset{p=1}{\overset{j}{}}}(L_k^{2𝔞}{\displaystyle \frac{k+2𝔞}{k}}L_{k1}^{2𝔞},L_p^{2𝔞}L_p^{2𝔞}(0))`$
$`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(k+2𝔞+1)}{k!}}\delta _{k1j}.`$
This completes the proof. ∎
Here we prove the following integral formula.
###### Lemma B.2
For $`j,\mathrm{}_0`$
$`{\displaystyle _0^z}𝑑xe^{x/2}x^𝔞G_j(x)=2e^{z/2}z^𝔞W_j(x),`$ (B.14)
$`{\displaystyle _0^z}𝑑xe^{x/2}x^𝔞F_2\mathrm{}(x)=2^{𝔞+1}\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{}+𝔞}{\mathrm{}}}\right)\gamma (𝔞+1,z/2)`$
$`2e^{z/2}z^𝔞\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{}+𝔞}{\mathrm{}}}\right){\displaystyle \underset{r=0}{\overset{\mathrm{}1}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{}r+𝔞}{\mathrm{}r}}\right)^1W_{2\mathrm{}2r}(z),`$ (B.15)
where $`\gamma (c,y),c>0`$ is the incomplete gamma function $`\gamma (c,y)=_0^y𝑑xe^xx^{c1}.`$
Remark 7. If we set $`𝔞=0`$ in (B.14), we will have the simpler equation
$$_0^z𝑑xe^{x/2}\frac{d}{dx}\left\{L_{j+1}^0(x)L_{j1}^0(x)\right\}=2e^{z/2}\left\{L_j^0(z)L_{j1}^0(z)\right\},j.$$
Proof of Lemma B.2. We first introduce the functions defined by
$`\psi _j^{2𝔞}(z)`$ $`=`$ $`{\displaystyle _0^z}𝑑xe^{x/2}x^𝔞{\displaystyle \frac{d}{dx}}L_j^{2𝔞}(x),j_0,`$
$`\phi _j^{2𝔞}(z)`$ $`=`$ $`{\displaystyle _0^z}𝑑xe^{x/2}x^𝔞L_j^{2𝔞}(x),j_0.`$
Then (B.3) gives
$$\phi _j^{2𝔞}(z)=\psi _{j+1}^{2𝔞}(z)+\psi _j^{2𝔞}(z).$$
(B.16)
On the other hand, by (B.2),
$`\psi _j^{2𝔞}(z)`$ $`=`$ $`j{\displaystyle _0^z}𝑑xe^{x/2}x^{𝔞1}W_j(x).`$
Noting the assumption $`𝔞>1`$ and the fact that $`W_j(x)=𝒪(x)`$, in $`x0`$,
$`2𝔞\psi _j^{2𝔞}(z)=2j\left\{e^{z/2}z^𝔞W_j(z)+{\displaystyle _0^z}𝑑xx^𝔞\left({\displaystyle \frac{1}{2}}e^{x/2}W_j(x)e^{x/2}{\displaystyle \frac{d}{dx}}W_j(x)\right)\right\}`$
$`=2je^{z/2}z^𝔞W_j(z)+j\phi _j^{2𝔞}(z)(j+2𝔞)\phi _{j1}^{2𝔞}(z)2j\psi _j^{2𝔞}(z)+2(j+2𝔞)\psi _{j1}^{2𝔞}(z).`$
Then we use (B.16) to eliminate $`\phi _j^{2𝔞}(z)`$ and $`\phi _{j1}^{2𝔞}(z)`$, and the equality (B.14) is obtained from the relation
$$_0^z𝑑xe^{x/2}x^𝔞G_j(x)=\psi _{j+1}^{2𝔞}(z)\frac{j+2𝔞}{j}\psi _{j1}^{2𝔞}(z).$$
Note that (B.14) gives a recurrence relation for $`\psi _j^{2𝔞}(z)`$,
$$\psi _{j+1}^{2𝔞}(z)=\frac{j+2𝔞}{j}\psi _{j1}^{2𝔞}(z)+2e^{z/2}z^𝔞W_j(z).$$
It is solved as
$`\psi _{2\mathrm{}+1}^{2𝔞}(z)`$ $`=`$ $`\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{}+𝔞}{\mathrm{}}}\right)\psi _1^{2𝔞}(z)+2e^{z/2}z^𝔞\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{}+𝔞}{\mathrm{}}}\right){\displaystyle \underset{r=0}{\overset{\mathrm{}1}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{}r+𝔞}{\mathrm{}r}}\right)^1W_{2\mathrm{}2r}(z).`$
Since $`L_1^{2𝔞}(x)=(1+2𝔞)x`$,
$`\psi _1^{2𝔞}(z)`$ $`=`$ $`{\displaystyle _0^z}𝑑xe^{x/2}x^𝔞{\displaystyle \frac{d}{dx}}L_1^{2𝔞}(z)={\displaystyle _0^z}𝑑xe^{x/2}x^𝔞=2^{𝔞+1}\gamma (𝔞+1,z/2).`$
Then we have (B.15). ∎
Then we can prove Lemma 4.2. From Lemma B.2 and two relations (4.15) and (B.7) we have
$`F_{2q},G_{2\mathrm{}+1}_{}`$ $`=`$ $`G_{2\mathrm{}+1},F_{2q}_{}=4(W_{2\mathrm{}+1},F_{2q}),`$
$`G_{2q+1},G_{2\mathrm{}+1}_{}`$ $`=`$ $`4(W_{2q+1},G_{2\mathrm{}+1})`$
$`=`$ $`4\left\{(W_{2q+1},F_{2\mathrm{}+1})+{\displaystyle \frac{2\mathrm{}+1+2𝔞}{2\mathrm{}+1}}(W_{2q+1},F_{2\mathrm{}1})\right\}.`$
Then (4.20) and (4.21) are derived from Lemma B.1. From Lemma B.2 and (B.9),
$`{\displaystyle _0^{\mathrm{}}}𝑑we^{w/2}w^𝔞\gamma (𝔞+1,w/2)F_2\mathrm{}(w)`$
$`=\left[\gamma (𝔞+1,w/2){\displaystyle _0^w}𝑑ze^{z/2}z^𝔞F_2\mathrm{}(z)\right]_0^{\mathrm{}}`$
$`{\displaystyle _0^{\mathrm{}}}𝑑w2^{(1+𝔞)}e^{w/2}w^𝔞{\displaystyle _0^w}𝑑ze^{z/2}z^𝔞F_2\mathrm{}(z)`$
$`=2^{𝔞+1}\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{}+𝔞}{\mathrm{}}}\right)\left\{\mathrm{\Gamma }(𝔞+1)^2{\displaystyle _0^{\mathrm{}}}𝑑w2^{(1+𝔞)}e^{w/2}w^𝔞\gamma (𝔞+1,w/2)\right\}`$
$`=2^{𝔞+1}\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{}+𝔞}{\mathrm{}}}\right)\left\{\mathrm{\Gamma }(𝔞+1)^2{\displaystyle \frac{1}{2}}\left[\gamma (𝔞+1,w/2)^2\right]_0^{\mathrm{}}\right\}=2^𝔞\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{}+𝔞}{\mathrm{}}}\right)\mathrm{\Gamma }(𝔞+1)^2.`$
Then by using Lemma B.2 and (B.9) we have (4.20), since
$`F_{2q},F_2\mathrm{}_{}`$ $`=`$ $`2{\displaystyle _0^{\mathrm{}}}𝑑we^{w/2}w^𝔞F_2\mathrm{}(w){\displaystyle _0^w}𝑑ze^{z/2}z^𝔞F_{2q}(z)`$
$`{\displaystyle _0^{\mathrm{}}}𝑑we^{w/2}w^𝔞F_2\mathrm{}(w){\displaystyle _0^{\mathrm{}}}𝑑ze^{z/2}z^𝔞F_{2q}(z)`$
$`=`$ $`2^{𝔞+2}\left({\displaystyle \genfrac{}{}{0pt}{}{q+𝔞}{q}}\right){\displaystyle _0^{\mathrm{}}}𝑑we^{w/2}w^𝔞\gamma (𝔞+1,w/2)F_2\mathrm{}(w)`$
$`2^{2𝔞+2}\left({\displaystyle \genfrac{}{}{0pt}{}{q+𝔞}{q}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{}+𝔞}{\mathrm{}}}\right)\mathrm{\Gamma }(𝔞+1)^2=0.`$
This completes the proof. ∎
## Appendix C Inverse of $`\{\alpha _{k,j}\}`$
We put
$$Q_2\mathrm{}(x)=F_2\mathrm{}(x)\text{and}Q_{2\mathrm{}+1}(x)=G_{2\mathrm{}+1}(x),$$
(C.1)
for $`\mathrm{}_0`$, and $`Q_k0`$ for $`k_{}`$. Lemma 4.1 gives the expansion formula of $`Q_k(x)`$ in terms of $`\{L_j^\nu (x)\}`$. Here we give the formula to expand $`L_j^\nu (x)`$ in terms of $`\{Q_k(x)\}`$. In other words, we calculate the inverse of the matrix $`𝜶=(\alpha _{k,j})`$ given by (4.16), which is denoted by $`𝜷=(\beta _{j,k})`$ and used in Section 5.2.
Let $`b(n)=(n+2𝔞)/n,n,`$ and
$$b(m,n)=\{\begin{array}{cc}b(m)b(m+2)\mathrm{}b(n),\hfill & \text{if }m,n\text{ are odd and }mn,\hfill \\ 1,\hfill & \text{if }m,n\text{ are odd and }m>n,\hfill \\ 0,\hfill & \text{otherwise}.\hfill \end{array}$$
(C.2)
Then the following lemma holds.
###### Lemma C.1
For $`j_0`$
$$L_j^\nu (x)=\underset{k=0}{\overset{j}{}}\beta _{j,k}Q_k(x).$$
(C.3)
where $`\beta _{j,k}`$, $`k,1,\mathrm{},j`$, $`j_0`$ are defined by the following:
When $`k`$ is even
$$\beta _{j,k}=\{\begin{array}{cc}0,\hfill & \text{if}j<k,\hfill \\ \left(\genfrac{}{}{0pt}{}{jk𝔟2}{jk}\right),\hfill & \text{if}jk,\hfill \end{array}$$
(C.4)
and, when $`k`$ is odd
$$\beta _{j,k}=\{\begin{array}{cc}0,\hfill & \text{if}j<k,\hfill \\ \underset{r=[(k+1)/2]}{\overset{[(j+1)/2]}{}}b(k+2,2r1)\left(\genfrac{}{}{0pt}{}{j2r𝔟1}{j2r+1}\right),\hfill & \text{if}jk.\hfill \end{array}$$
(C.5)
Proof. From (B.5) and (B.6) we have
$`Q_2\mathrm{}(x)={\displaystyle \underset{j=0}{\overset{2\mathrm{}}{}}}L_j^{2𝔞}(x),Q_{2\mathrm{}+1}(x)={\displaystyle \underset{j=0}{\overset{2\mathrm{}+1}{}}}L_j^{2𝔞}(x)+b(2\mathrm{}+1){\displaystyle \underset{j=0}{\overset{2\mathrm{}1}{}}}L_j^{2𝔞}(x).`$
By simple calculations we see that
$$L_j^{2𝔞}(x)=\underset{k=0}{\overset{j}{}}\overline{\beta }_{j,k}Q_k(x),j_0$$
(C.6)
where $`\overline{\beta }_{j,k}`$, $`k=0,1,\mathrm{},j`$, $`j_0`$, are defined by the following:
When $`k`$ is even
$$\overline{\beta }_{j,k}=\{\begin{array}{cc}1,\hfill & \text{if}j=k,\hfill \\ 1\hfill & \text{if}j=k+1,\hfill \\ 0,\hfill & \text{otherwise},\hfill \end{array}$$
and, when $`k`$ is odd
$$\overline{\beta }_{j,k}=\{\begin{array}{cc}b(k+2,j1),\hfill & \text{if }j\text{ is even},\hfill \\ b(k+2,j),\hfill & \text{if }j\text{ is odd}.\hfill \end{array}$$
Using the formula (B.4), (C.6) gives
$`L_j^\nu (x)`$ $`=`$ $`{\displaystyle \underset{p=0}{\overset{j}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{jp𝔟1}{jp}}\right){\displaystyle \underset{k=0}{\overset{p}{}}}\overline{\beta }_{p,k}Q_k(x).`$
If we define
$$\beta _{j,k}=\underset{p=k}{\overset{j}{}}\overline{\beta }_{p,k}\left(\genfrac{}{}{0pt}{}{jp𝔟1}{jp}\right)$$
(C.7)
for $`jk`$ and $`\beta _{j,k}=0`$ for $`j<k`$, (C.3) is satisfied. The expressions (C.4) and (C.5) are derived from (C.7) by simple calculation. ∎
## Appendix D Elementary Calculation for Asymptotics Estimation
By Stirling’s formula $`\mathrm{\Gamma }(x)\sqrt{2\pi }x^{x1/2}e^x,x\mathrm{}`$, we have
$`{\displaystyle \frac{\mathrm{\Gamma }(n+\alpha +1)}{\mathrm{\Gamma }(n+1)}}`$ $``$ $`(n+1)^\alpha ,n\mathrm{},`$ (D.1)
$`\left({\displaystyle \genfrac{}{}{0pt}{}{n+\alpha }{n}}\right)`$ $``$ $`{\displaystyle \frac{(n+1)^\alpha }{\mathrm{\Gamma }(\alpha +1)}}n\mathrm{},`$ (D.2)
for any $`\alpha _{}`$, and
$`b(1,2\mathrm{}1)={\displaystyle \frac{\mathrm{\Gamma }(2𝔞+1)}{2^𝔞\mathrm{\Gamma }(𝔞+1)}}\left({\displaystyle \genfrac{}{}{0pt}{}{2\mathrm{}+2𝔞}{2\mathrm{}}}\right)/\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{}+𝔞}{\mathrm{}}}\right)(2\mathrm{})^𝔞,\mathrm{}\mathrm{}`$ (D.3)
$`b(2\mathrm{}+1,2p1)={\displaystyle \frac{b(1,2p1)}{b(1,2\mathrm{}1)}}\left({\displaystyle \frac{p}{\mathrm{}}}\right)^𝔞,\mathrm{}\mathrm{}`$ (D.4)
for $`\mathrm{},p`$ with $`\mathrm{}<p`$, where $`b(m,n)`$ is defined by (C.2).
From now on, we assume that $`T=N`$, $`t_m=T+s_m`$ with $`s_m<0`$. We set
$$2\mathrm{}=N\theta \text{and}j=2\mathrm{}\eta ,$$
(D.5)
and consider the limit $`N\mathrm{}`$ with some $`\eta ,\theta (0,\mathrm{})`$. Then we have
$$\chi _m^j=\left(\frac{2Tt_m}{t_m}\right)^j=\left(1\frac{2s_m}{t_m}\right)^{N\theta \eta }\mathrm{exp}\left(2s_m\theta \eta \right),N\mathrm{},$$
and
$$\underset{p=0}{\overset{\alpha }{}}(1)^p\left(\genfrac{}{}{0pt}{}{\alpha }{p}\right)\chi _m^{jp}(2\mathrm{})^\alpha \left(\frac{d}{d\eta }\right)^\alpha \mathrm{exp}\left(2s_m\theta \eta \right),N\mathrm{}.$$
(D.6)
We use the following identities (see (4.10) and pages 8, 201 and 202 in ).
(i) Let $`\alpha _0`$ and $`c`$. Then
$$\underset{p=0}{\overset{\alpha }{}}(1)^p\left(\genfrac{}{}{0pt}{}{\alpha }{p}\right)\left(\genfrac{}{}{0pt}{}{np+c}{npj}\right)=\left(\genfrac{}{}{0pt}{}{n\alpha +c}{nj}\right).$$
(D.7)
(ii) Let $`\alpha _0`$, $`c`$ and $`a_k`$, $`k=1,2,\mathrm{},`$ be a sequence in $``$. Then
$`{\displaystyle \underset{r=0}{\overset{\mathrm{}}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{r+c}{r}}\right)a_r={\displaystyle \underset{r=0}{\overset{\mathrm{}}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{r+c+\alpha }{r}}\right){\displaystyle \underset{p=0}{\overset{\alpha }{}}}(1)^p\left({\displaystyle \genfrac{}{}{0pt}{}{\alpha }{p}}\right)a_{r+p},`$ (D.8)
(iii) Let $`\alpha _0`$, and $`a_k,b_k`$, $`k=1,2,\mathrm{},`$ be sequences in $``$. Then
$`{\displaystyle \underset{k=0}{\overset{\alpha }{}}}(1)^k\left({\displaystyle \genfrac{}{}{0pt}{}{\alpha }{k}}\right)a_kb_k={\displaystyle \underset{\beta =0}{\overset{\alpha }{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{\alpha }{\beta }}\right){\displaystyle \underset{p=0}{\overset{\beta }{}}}(1)^p\left({\displaystyle \genfrac{}{}{0pt}{}{\beta }{p}}\right)a_p{\displaystyle \underset{q=0}{\overset{\alpha \beta }{}}}(1)^q\left({\displaystyle \genfrac{}{}{0pt}{}{\alpha \beta }{q}}\right)b_q.`$ (D.9)
###### Lemma D.1
For any $`\alpha _0`$ and $`w0`$ we have
$$\underset{p=0}{\overset{\alpha }{}}(1)^p\left(\genfrac{}{}{0pt}{}{\alpha }{p}\right)L_{np}^\nu \left(\frac{w}{n}\right)\left(\frac{w}{n}\right)^{\alpha \nu }\left(\frac{d}{dw}\right)^\alpha \left\{w^{\nu /2}J_\nu (2\sqrt{w})\right\},n\mathrm{}.$$
Proof. From the definition of the Laguerre polynomials (4.12) and (D.7), we have
$`{\displaystyle \underset{p=0}{\overset{\alpha }{}}}(1)^p\left({\displaystyle \genfrac{}{}{0pt}{}{\alpha }{p}}\right)L_{np}^\nu \left(y\right)`$ $`=`$ $`{\displaystyle \underset{p=0}{\overset{\alpha }{}}}(1)^p\left({\displaystyle \genfrac{}{}{0pt}{}{\alpha }{p}}\right){\displaystyle \underset{j=0}{\overset{np}{}}}(1)^j\left({\displaystyle \genfrac{}{}{0pt}{}{np+\nu }{npj}}\right){\displaystyle \frac{y^j}{j!}}`$
$`=`$ $`{\displaystyle \underset{j=0}{\overset{n}{}}}{\displaystyle \frac{(1)^j}{j!}}\left({\displaystyle \genfrac{}{}{0pt}{}{n\alpha +\nu }{nj}}\right)y^j.`$
Hence, by (D.2)
$`\underset{n\mathrm{}}{lim}\left({\displaystyle \frac{w}{n}}\right)^{\nu \alpha }{\displaystyle \underset{p=0}{\overset{\alpha }{}}}(1)^p\left({\displaystyle \genfrac{}{}{0pt}{}{\alpha }{p}}\right)L_{np}^\nu \left({\displaystyle \frac{w}{n}}\right)`$ $`=`$ $`\underset{n\mathrm{}}{lim}{\displaystyle \underset{j=0}{\overset{n}{}}}{\displaystyle \frac{(1)^j}{j!}}\left({\displaystyle \genfrac{}{}{0pt}{}{n\alpha +\nu }{nj}}\right)\left({\displaystyle \frac{w}{n}}\right)^{j+\nu \alpha }`$
$`=`$ $`{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}(1)^j{\displaystyle \frac{w^{\nu +j\alpha }}{\mathrm{\Gamma }(j\alpha +\nu +1)j!}}`$
$`=`$ $`\left({\displaystyle \frac{d}{dw}}\right)^\alpha \left({\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}(1)^j{\displaystyle \frac{w^{\nu +j}}{\mathrm{\Gamma }(\nu +j+1)j!}}\right)`$
Then we obtain the lemma ∎
Applying the above lemma, we obtain the following asymptotics, where $`L_j^\nu `$ and $`\widehat{L}_j^\nu `$ are defined by (5.1).
###### Lemma D.2
For any $`\alpha _0`$, $`\theta ,\eta (0,\mathrm{})`$, and $`x_+`$ we have
$`{\displaystyle \underset{p=0}{\overset{\alpha }{}}}(1)^p\left({\displaystyle \genfrac{}{}{0pt}{}{\alpha }{p}}\right)L_{jp}^\nu ({\displaystyle \frac{x}{N}},s_m){\displaystyle \frac{(2\mathrm{})^{\nu \alpha }}{(\theta x)^\nu }}\stackrel{~}{J}_\nu ^{(\alpha )}(\theta ,\eta ,x,s_m),N\mathrm{},`$ (D.10)
$`{\displaystyle \underset{p=0}{\overset{\alpha }{}}}(1)^p\left({\displaystyle \genfrac{}{}{0pt}{}{\alpha }{p}}\right)\widehat{L}_{j+p}^\nu ({\displaystyle \frac{x}{N}},s_m)(2\mathrm{})^\alpha \widehat{J}_\nu ^{(\alpha )}(\theta ,\eta ,x,s_m),N\mathrm{}.`$ (D.11)
Proof. Since
$`{\displaystyle \underset{p=0}{\overset{\alpha }{}}}(1)^p\left({\displaystyle \genfrac{}{}{0pt}{}{\alpha }{p}}\right)L_{jp}^\nu \left({\displaystyle \frac{x}{N}}\right)\chi _m^{jp}`$
$`={\displaystyle \underset{\beta =0}{\overset{\alpha }{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{\alpha }{\beta }}\right){\displaystyle \underset{p=0}{\overset{\beta }{}}}(1)^p\left({\displaystyle \genfrac{}{}{0pt}{}{\beta }{p}}\right)L_{jp}^\nu \left({\displaystyle \frac{x}{N}}\right){\displaystyle \underset{q=0}{\overset{\alpha \beta }{}}}(1)^q\left({\displaystyle \genfrac{}{}{0pt}{}{\alpha \beta }{q}}\right)\chi _m^{jq},`$
by (D.9), the asymptotic (D.10) is derived from (D.6) and Lemma D.1 with $`n=j=N\theta \eta ,w=\theta \eta x`$. From (D.9), we have
$`{\displaystyle \underset{k=0}{\overset{\alpha }{}}}(1)^k\left({\displaystyle \genfrac{}{}{0pt}{}{\alpha }{k}}\right)\widehat{L}_{j+k}^\nu ({\displaystyle \frac{x}{N}},s_m)`$
$`={\displaystyle \underset{\beta =0}{\overset{\alpha }{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{\alpha }{\beta }}\right){\displaystyle \underset{p=0}{\overset{\beta }{}}}(1)^p\left({\displaystyle \genfrac{}{}{0pt}{}{\beta }{p}}\right)L_{j+p}^\nu \left({\displaystyle \frac{x}{N}}\right)\chi _m^{(j+p)}{\displaystyle \underset{q=0}{\overset{\alpha \beta }{}}}(1)^q\left({\displaystyle \genfrac{}{}{0pt}{}{\alpha \beta }{q}}\right){\displaystyle \frac{\mathrm{\Gamma }(j+q+1)}{\mathrm{\Gamma }(j+q+1+\nu )}}.`$
(D.12)
By (D.1), we see
$`{\displaystyle \underset{q=0}{\overset{\alpha \beta }{}}}(1)^q\left({\displaystyle \genfrac{}{}{0pt}{}{\alpha \beta }{q}}\right){\displaystyle \frac{\mathrm{\Gamma }(j+q+1)}{\mathrm{\Gamma }(j+q+1+\nu )}}`$ $`=`$ $`\nu {\displaystyle \underset{q=0}{\overset{\alpha \beta 1}{}}}(1)^q\left({\displaystyle \genfrac{}{}{0pt}{}{\alpha \beta 1}{q}}\right){\displaystyle \frac{\mathrm{\Gamma }(j+q+1)}{\mathrm{\Gamma }(j+q+2+\nu )}}`$
$`=`$ $`{\displaystyle \frac{\nu (\nu +1)\mathrm{}(\nu +\alpha \beta 1)\mathrm{\Gamma }(j+1)}{\mathrm{\Gamma }(j+1+\alpha \beta +\nu )}},`$
$``$ $`(2\mathrm{})^{(\alpha \beta +\nu )}\left({\displaystyle \frac{d}{d\eta }}\right)^{\alpha \beta }\eta ^\nu ,N\mathrm{}.`$
On the other hand, (D.10) gives
$`{\displaystyle \underset{p=0}{\overset{\beta }{}}}(1)^p\left({\displaystyle \genfrac{}{}{0pt}{}{\beta }{p}}\right)L_{j+p}^\nu \left({\displaystyle \frac{x}{N}}\right)\chi _m^{(j+p)}(2\mathrm{})^{\nu \beta }\left({\displaystyle \frac{d}{d\eta }}\right)^\beta \left\{\eta ^\nu \widehat{J}_\nu (\theta ,\eta ,x,s_m)\right\}.`$
Hence, the asymptotic (D.11) is derived from (D.12). ∎
## Appendix E On Temporally Homogeneous Limit
###### Lemma E.1
For any $`c`$ and $`\eta ,\theta ,x0`$, we have that as $`t\mathrm{}`$
$`\stackrel{~}{J}_\nu ^{(c)}(\theta ,\eta ,x,t)(2t\theta )^c(\theta \eta x)^{\nu /2}J_\nu (2\sqrt{\theta \eta x})e^{2t\theta \eta },`$ (E.1)
$`\widehat{J}_\nu ^{(c)}(\theta ,\eta ,x,t)(2t\theta )^c(\theta \eta x)^{\nu /2}J_\nu (2\sqrt{\theta \eta x})e^{2t\theta \eta },`$ (E.2)
$`{\displaystyle _1^{\mathrm{}}}𝑑\xi \xi ^𝔞\widehat{J}_\nu ^{(c+1)}(\theta ,\xi ,x,t)(2t\theta )^c(\theta x)^{\nu /2}J_\nu (2\sqrt{\theta x})e^{2t\theta }.`$ (E.3)
Proof. From the expression (2.16) with the definition (2.12), we have
$`\stackrel{~}{J}_\nu ^{(c)}(\theta ,\eta ,x,t)={\displaystyle \frac{e^{2t\theta \eta }}{\mathrm{\Gamma }(c)}}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^k\eta ^{kc}}{k!(kc)}}{\displaystyle \underset{j=0}{\overset{k}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{k}{j}}\right)\stackrel{~}{J}_\nu ^{(j)}(\theta ,\eta ,x,0)(2t\theta )^{kj}`$
$`{\displaystyle \frac{e^{2t\theta \eta }}{\mathrm{\Gamma }(c)}}(2t\theta )^c\stackrel{~}{J}_\nu (\theta ,\eta ,x,0){\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^k(2t\theta \eta )^{kc}}{k!(kc)}},t\mathrm{}.`$
From the relation
$$\frac{d}{dz}\underset{k=0}{\overset{\mathrm{}}{}}\frac{(1)^kz^{kc}}{k!(kc)}=z^{c1}e^z,$$
and the equation
$$\mathrm{\Gamma }(c)=\underset{k=0}{\overset{\mathrm{}}{}}\frac{(1)^k}{k!(kc)}+_1^{\mathrm{}}𝑑zz^{c1}e^z$$
(see (1.1.19) in ), we have
$$\mathrm{\Gamma }(c)=\underset{t\mathrm{}}{lim}\underset{k=0}{\overset{\mathrm{}}{}}\frac{(1)^k(2t\theta \eta )^{kc}}{k!(kc)}.$$
Then we conclude
$`\stackrel{~}{J}_\nu ^{(c)}(\theta ,\eta ,x,t)e^{2t\theta \eta }(2t\theta )^c\stackrel{~}{J}_\nu (\theta ,\eta ,x,0)=(2t\theta )^c\stackrel{~}{J}_\nu (\theta ,\eta ,x,t),t\mathrm{}.`$
Hence (E.1) is derived from (2.12).
Let $`n=[c+1]_+`$ and $`\beta >0`$ with $`c=n\beta `$. Since
$$\left(\frac{d}{d\eta }\right)^n\widehat{J}_\nu (\theta ,\eta ,x,t)(2t\theta )^n\widehat{J}_\nu (\theta ,\eta ,x,t),t\mathrm{},$$
(2.17) gives
$`\widehat{J}_\nu ^{(c)}(\theta ,\eta ,x,t){\displaystyle \frac{(2t\theta )^n}{\mathrm{\Gamma }(\beta )}}{\displaystyle _0^{\mathrm{}}}𝑑\xi \xi ^{\beta 1}\widehat{J}_\nu (\theta ,\eta +\xi ,x,t)`$
$`={\displaystyle \frac{(2t\theta )^{n\beta }}{\mathrm{\Gamma }(\beta )}}{\displaystyle _0^{\mathrm{}}}𝑑\zeta \zeta ^{\beta 1}\widehat{J}_\nu (\theta ,\eta +{\displaystyle \frac{\zeta }{2t\theta }},x,t)`$
$`(2t\theta )^c\widehat{J}_\nu (\theta ,\eta ,x,t),t\mathrm{}.`$
Then (E.2) is derived from (2.13).
From (E.2) we have
$`{\displaystyle _1^{\mathrm{}}}𝑑\xi \xi ^𝔞\widehat{J}_\nu ^{(c+1)}(\theta ,\xi ,x,t)`$ $``$ $`(2t\theta )^{c+1}{\displaystyle _1^{\mathrm{}}}𝑑\xi \xi ^𝔞\widehat{J}_\nu (\theta ,\xi ,x,0)e^{2t\theta \xi }`$
$``$ $`(2t\theta )^c\widehat{J}_\nu (\theta ,1,x,t),t\mathrm{}.`$
This completes the proof. ∎
Applying the above lemma, we have as $`s_m,s_n\mathrm{}`$ with the difference $`s_ns_m`$ fixed
$`𝒟(s_m,x;s_n,y)`$ $``$ $`{\displaystyle \frac{(xy)^{𝔟/2}(s_ns_m)}{2^{2𝔟+3}(s_ms_n)^{𝔟+1}}}{\displaystyle _0^1}𝑑\theta \theta ^𝔟J_\nu (2\sqrt{\theta x})J_\nu (2\sqrt{\theta y})e^{2(s_m+s_n)\theta }`$
$``$ $`{\displaystyle \frac{(xy)^{𝔟/2}(s_ms_n)}{2^{2𝔟+4}(s_m+s_n)(s_ms_n)^{𝔟+1}}}J_\nu (2\sqrt{x})J_\nu (2\sqrt{y})e^{2(s_m+s_n)},`$
$`\stackrel{~}{}(s_m,x;s_n,y)`$ $``$ $`{\displaystyle \frac{2^{2𝔟+1}(s_ms_n)^𝔟(s_ns_m)}{(xy)^{𝔟/2}}}{\displaystyle _1^{\mathrm{}}}𝑑\theta \theta ^𝔟J_\nu (2\sqrt{\theta x})J_\nu (2\sqrt{\theta y})e^{2(s_m+s_n)\theta }`$
$``$ $`{\displaystyle \frac{2^{2𝔟}(s_ms_n)^𝔟(s_ns_m)}{(s_m+s_n)(xy)^{𝔟/2}}}J_\nu (2\sqrt{x})J_\nu (2\sqrt{y})e^{2(s_m+s_n)},`$
and
$`𝒮(s_m,x;s_n,y)\left({\displaystyle \frac{y}{x}}\right)^{𝔟/2}{\displaystyle _0^1}𝑑\theta J_\nu (2\sqrt{\theta x})J_\nu (2\sqrt{\theta y})e^{2(s_ms_n)\theta }.`$
It is then clear that
$$\underset{s_m,s_n\mathrm{}}{lim}𝒟(s_m,x;s_n,y)\stackrel{~}{}(s_m,x;s_n,y)=0.$$
## Acknowledgment
The authors would like to thank T. Sasamoto and T. Imamura for useful discussion on determinantal and Pfaffian processes.
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# Schur-Weyl reciprocity between the quantum superalgebra and the Iwahori-Hecke algebra
## 1 Introduction
In the representation theory, the classification and the construction of the irreducible representations are essential themes. In the first half of the twentieth century, I. Schur introduced a prominent method to obtain the finite dimensional irreducible representations of the general linear group $`\mathrm{GL}(n,)`$, or equivalently of its Lie algebra $`𝔤𝔩(n,)`$, which we call Schur-Weyl reciprocity at present. Schur applied this method to the permutation action of the symmetric group $`𝔖_r`$ and the diagonal action of $`\mathrm{GL}(n,)`$ on the tensor powers $`V^r`$ of the $`n`$ dimensional complex vector space $`V`$.
After this work, Schur-Weyl reciprocity has been extended to various groups and algebras. Brauer obtained the centralizer algebra of the orthogonal Lie group $`O(n)`$. Sergeev and Berele-Regev extended the Schur’s result to the general super Lie algebra $`𝔤𝔩(m,n)`$. Jimbo extended it to the $`q`$-analogue case. He established Schur-Weyl reciprocity between the quantum enveloping algebra $`U_q(𝔤𝔩_{n+1})`$ and the Iwahori-Hecke algebra of type $`A`$. As in the book of Curtis-Reiner, the representation theory of Iwahori-Hecke algebras is an important part in representation theories of finite groups of Lie type. Hence we will focus on the representation theory of Iwahori-Hecke algebras.
In , we defined a $`q`$-deformation of the alternating group as a subalgebra of the Iwahori-Hecke algebra, and determined all the isomorphism classes of (ordinary) irreducible representations. After , we intended to compute character values of irreducible representations directly using combinatorial method. But this strategy did not go well because the notion of conjugacy classes of the $`q`$-deformation of the alternating group is obscure, hence we could not apply the classical($`q=1`$) case which is found in . Thereby we will take a representational approach to obtain the character table.
In , Regev obtained double centralizer properties for alternating groups. Roughly speaking, his result claims when one restrict the representation of the symmetric group on tensor space to the alternating group, the corresponding centralizer algebra enlarges in “super case” while does not change in “normal case”. Those facts suggest that Schur-Weyl reciprocity for $`𝔤𝔩(m,n)`$ is more suitable to describe the representation theory of the alternating group than that for $`𝔤𝔩(n)`$.
In this paper, we establish Schur-Weyl reciprocity between the quantum superalgebra $`U_q^\sigma \left(𝔤𝔩(m,n)\right)`$ and the Iwahori-Hecke algebra $`_{(q),r}(q)`$. We define the sign $`q`$-permutation representation of $`_{(q),r}(q)`$ on $`V^r`$ using an operator $`T`$ on $`V^2`$ defined by:
$$v_kv_lT=\{\begin{array}{cc}\frac{(1)^{|v_k|}(q+q^1)+qq^1}{2}v_kv_l\hfill & \text{if }k=l\text{,}\hfill \\ (1)^{|v_k||v_l|}v_lv_k+(qq^1)v_kv_l\hfill & \text{if }k<l\text{,}\hfill \\ (1)^{|v_k||v_l|}v_lv_k\hfill & \text{if }k>l\text{,}\hfill \end{array}$$
(1.1)
where $`V`$ is an $`(m+n)`$ dimensional $`_2`$-graded $`(q)`$-vector space and $`||`$ is the degree map. This action reduces to the sign permutation action(see ) of the symmetric group when $`q1`$ and to the well-known action of $`_{(q),r}(q)`$ obtained from Drinfeld-Jimbo solutions to the Yang-Baxter equation when $`n=0`$.
The quantum superalgebra has been defined in several articles such as , or . $`U_q^\sigma \left(𝔤𝔩(m,n)\right)`$ is a Hopf algebra obtained from the “naive” quantum superalgebra $`U_q\left(𝔤𝔩(m,n)\right)`$, which is a Hopf superalgebra, by adding an involutive element $`\sigma `$. We show that the vector representation of $`U_q^\sigma \left(𝔤𝔩(m,n)\right)`$ on $`V^r`$, which is found in , and the sign $`q`$-permutation representation of $`_{(q),r}(q)`$ are commutants of one another in $`V^r`$. Furthermore, extending the base field to the algebraic closure and applying the double centralizer theorem, we obtain the tensor space decomposition of $`\left(V\overline{(q)}\right)^r`$ as $`_{\overline{(q)},r}(q)\left(U_q^\sigma \left(𝔤𝔩(m,n)\right)\overline{(q)}\right)`$-modules. Our result will be very useful to the representation theory of the $`q`$-deformation of the alternating group.
## 2 The sign $`q`$-permutation representation of the Iwahori-Hecke algebra of type $`A`$
In this section, we shall define the sign $`q`$-permutation representation of the Iwahori-Hecke algebra of type $`A`$, which is a $`q`$-deformation of the sign permutation module introduced by several precedent works such as , .
Let $`(W,S=\{s_1,\mathrm{},s_r\})`$ be a Coxeter system of rank $`r`$. Let $`R`$ be a commutative domain with $`1`$, and let $`q_i(i=1,\mathrm{},r)`$ be any invertible elements of $`R`$ such that $`q_i=q_j`$ if $`s_i`$ is conjugate to $`s_j`$ in $`W`$. The Iwahori-Hecke algebra $`_R(W,S)`$ is an $`R`$-algebra generated by $`\{T_{s_i}|s_iS\}`$ with the relations:
1. $`T_{s_i}^2=(q_iq_i^1)T_{s_i}+1`$ if $`i=1,2,\mathrm{},r`$,
2. $`(T_{s_i}T_{s_j})^{k_{ij}}=(T_{s_j}T_{s_i})^{k_{ij}}`$ if $`m_{ij}=2k_{ij}`$,
3. $`(T_{s_i}T_{s_j})^{k_{ij}}T_{s_i}=(T_{s_j}T_{s_i})^{k_{ij}}T_{s_j}`$ if $`m_{ij}=2k_{ij}+1`$,
where $`m_{ij}`$ is the order of $`s_is_j`$ in $`W`$. We define $`T_w=T_{s_{i_1}}T_{s_{i_2}}\mathrm{}T_{s_{i_k}}`$ where $`w=s_{i_1}s_{i_2}\mathrm{}s_{i_k}`$ is a reduced expression of $`w`$. It is known that $`T_w`$ is well defined because two elements $`T_w`$ and $`T_w^{}`$, where $`w`$ and $`w^{}`$ are reduced expressions of an element of $`W`$, coincide and that $`\{T_w|wW\}`$ form a basis of $`_R(W,S)`$ as free $`R`$-modules. The relations (H1)–(H3) is equivalent to the following two relations:
1. $`T_{s_i}T_w=T_{s_iw}`$ if $`l(w)<l(s_iw)`$,
2. $`T_{s_i}T_w=(q_iq_i^1)T_w+T_{s_iw}`$ if $`l(w)>l(s_iw)`$,
or equivalently,
1. $`T_wT_{s_i}=T_{ws_i}`$ $`ifl(w)<l(ws_i)`$,
2. $`T_wT_{s_i}=(q_iq_i^1)T_w+T_{ws_i}`$ if $`l(w)>l(ws_i)`$,
where $`l(w)`$ means the length of $`w`$. We write $`T_i=T_{s_i}`$ for brevity.
If $`(W,S)`$ is of type $`A`$ and of rank $`r1`$, then $`W`$ is isomorphic to the symmetric group $`𝔖_r`$. Furthermore, all the elements of $`S`$ are conjugate to each other, hence we may assume $`q_1=\mathrm{}=q_{r1}=q`$. The Iwahori-Hecke algebra $`_{R,r}(q)=_R(W,S)`$ of type $`A`$ has defining relations:
1. $`T_i^2=(qq^1)T_i+1`$ if $`i=1,2,\mathrm{},r1`$,
2. $`T_iT_{i+1}T_i=T_{i+1}T_iT_{i+1}`$ if $`i=1,2,\mathrm{},r2`$,
3. $`T_iT_j=T_jT_i`$ if $`|ij|>1`$.
Let $`V=_{k=1}^{m+n}Rv_k`$ be a $`_2`$-graded $`R`$-module of rank $`m+n`$. By $`_2`$-graded, we mean that $`V`$ is a direct sum of two submodules $`V_{\overline{0}}=_{k=1}^mRv_k`$ and $`V_{\overline{1}}=_{k=m+1}^{m+n}Rv_k`$, and that for each homogeneous element the degree map $`||`$
$$|v|=\{\begin{array}{cc}0\hfill & \text{if }vV_{\overline{0}},\hfill \\ 1\hfill & \text{if }vV_{\overline{1}}\text{,}\hfill \end{array}$$
is given. In order to define a representation of $`_{R,r}(q)`$ on $`V^r`$, we define a right operator $`T`$ on $`VV`$ as follows.
$$v_kv_lT=\{\begin{array}{cc}\frac{(1)^{|v_k|}(q+q^1)+qq^1}{2}v_kv_l\hfill & \text{if }k=l\text{,}\hfill \\ (1)^{|v_k||v_l|}v_lv_k+(qq^1)v_kv_l\hfill & \text{if }k<l\text{,}\hfill \\ (1)^{|v_k||v_l|}v_lv_k\hfill & \text{if }k>l\text{.}\hfill \end{array}$$
(2.1)
The factor $`2^1`$ in the case $`k=l`$ vanishes whether $`v_kV_{\overline{0}}`$ or $`v_kV_{\overline{1}}`$. So the above definition makes sense on $`R`$. Now we obtain right operators $`\mathrm{Id}^{i1}T\mathrm{Id}^{ri1}`$ ($`i=1,2,\mathrm{},r1`$) on $`V^r`$ where $`\mathrm{Id}`$ is the identity operator on $`V`$. Let us define a map $`\pi _r:\{T_1,\mathrm{},T_{r1}\}\mathrm{End}_R(V^r)`$ by $`\pi _r(T_i)=\mathrm{Id}^{i1}T\mathrm{Id}^{ri1}`$ ($`i=1,2,\mathrm{},r1`$).
###### Theorem 2.1.
$`\pi _r`$ defines a representation of $`_{R,r}(q)`$ on $`V^r`$.
###### Proof.
One can check that the above operators satisfy the defining relations (A1)–(A3) by a direct computation. For example, the relation (A1) is shown as follows.
case1 : $`k=l`$
$$\begin{array}{cc}\hfill v_kv_kT^2& =\frac{1}{4}\left\{(1)^{|v_k|}(q+q^1)+qq^1\right\}^2v_kv_k\hfill \\ & =\frac{1}{2}\left\{q^2+q^2+(1)^{|v_k|}(q+q^1)(qq^1)\right\}v_kv_k\hfill \end{array}$$
$$\begin{array}{cc}\hfill v_kv_k\left\{(qq^1)T+1\right\}& =\frac{(qq^1)}{2}\left\{(1)^{|v_k|}(q+q^1)+qq^1\right\}v_kv_k+v_kv_k\hfill \\ & =\frac{1}{2}\left\{q^2+q^2+(1)^{|v_k|}(q+q^1)(qq^1)\right\}v_kv_k\hfill \end{array}$$
case2 : $`k<l`$
$$\begin{array}{cc}\hfill v_kv_lT^2& =\left\{(1)^{|v_k||v_l|}v_lv_k+(qq^1)v_kv_l\right\}T\hfill \\ & =v_kv_l+(qq^1)\left\{(1)^{|v_k||v_l|}v_lv_k+(qq^1)v_kv_l\right\}\hfill \\ & =(qq^1)(1)^{|v_k||v_l|}v_lv_k+(q^2+q^21)v_kv_l\hfill \end{array}$$
$$\begin{array}{cc}\hfill v_kv_l\left\{(qq^1)T+1\right\}& =(qq^1)\left\{(1)^{|v_k||v_l|}v_lv_k+(qq^1)v_kv_l\right\}+v_kv_l\hfill \\ & =(qq^1)(1)^{|v_k||v_l|}v_lv_k+(q^2+q^21)v_kv_l\hfill \end{array}$$
case3 : $`k>l`$
$$\begin{array}{cc}\hfill v_kv_lT^2& =(1)^{|v_k||v_l|}\left\{(1)^{|v_k||v_l|}v_kv_l+(qq^1)v_lv_k\right\}\hfill \\ & =v_kv_l+(1)^{|v_k||v_l|}(qq^1)v_lv_k\hfill \end{array}$$
$$\begin{array}{cc}\hfill v_kv_l\left\{(qq^1)T+1\right\}& =(qq^1)(1)^{|v_k||v_l|}v_lv_k+v_kv_l\hfill \end{array}$$
(A2) can be shown in a similar manner to (A1), albeit slightly lengthy. The relation (A3) is obvious. ∎
###### Remark 2.2.
Another definition of the Iwahori-Hecke algebra, which is frequently used
1. $`\overline{T}_i^2=(q^{}1)\overline{T}_i+q^{}`$ if $`i=1,2,\mathrm{},r1`$,
2. $`\overline{T}_i\overline{T}_{i+1}\overline{T}_i=\overline{T}_{i+1}\overline{T}_i\overline{T}_{i+1}`$ if $`i=1,2,\mathrm{},r2`$,
3. $`\overline{T}_i\overline{T}_j=\overline{T}_j\overline{T}_i`$ if $`|ij|>1`$,
can be obtained from previous definition (A1)–(A3) by letting $`q^{}=q^2`$ and $`\overline{T}_i=qT_i`$. Accordingly, we get another right operator $`\overline{T}`$ on $`VV`$ as follows.
$$v_kv_l\overline{T}=\{\begin{array}{cc}\frac{(1)^{|v_k|}(q^{}+1)+q^{}1}{2}v_kv_k\hfill & \text{if }k=l\text{,}\hfill \\ (1)^{|v_k||v_l|}\sqrt{q^{}}v_lv_k+(q^{}1)v_kv_l\hfill & \text{if }k<l\text{,}\hfill \\ (1)^{|v_k||v_l|}\sqrt{q^{}}v_lv_k\hfill & \text{if }k>l\text{.}\hfill \end{array}$$
(2.2)
We can also define the action of the alternative generators $`\overline{T}_i`$ of the Iwahori-Hecke algebra in a similar manner.
This representation $`\pi _r`$ is reduced to the (normal) $`q`$-permutation representation of $`_{R,r}(q)`$ obtained from Drinfeld-Jimbo solutions to the Yang-Baxter equation when $`n=0`$ and to the sign permutation action (see ) of the symmetric group when $`q1`$.
## 3 The quantum superalgebra $`U_q^\sigma \left(𝔤𝔩(m,n)\right)`$ and the vector representation
In Kac’s paper, classical superalgebras have been classified and studied in detail. Quantum superalgebras have been defined in several articles such as , or . Each definition of them seems to be based on essentially. $`U_q^\sigma \left(𝔤𝔩(m,n)\right)`$ is a Hopf algebra obtained from the “naive” quantum superalgebra $`U_q\left(𝔤𝔩(m,n)\right)`$, which is a Hopf superalgebra, by adding an involutive element $`\sigma `$. According to , we adopt $`U_q^\sigma \left(𝔤𝔩(m,n)\right)`$ to construct the vector representation on the tensor space $`V^r`$.
Let $`\mathrm{\Pi }=\{\alpha _i\}_{iI}`$ be a set of simple roots with the index set $`I=\{1,\mathrm{},r\}`$. We assume that $`I`$ is a disjoint union of two subsets $`I_{\mathrm{even}}`$ and $`I_{\mathrm{odd}}`$. We define a map $`p:I\{0,1\}`$ to be such that
$$p(i)=\{\begin{array}{cc}0\hfill & \text{if }iI_{\mathrm{even}},\hfill \\ 1\hfill & \text{if }iI_{\mathrm{odd}}\text{.}\hfill \end{array}$$
Let $`P`$ be a free $``$-module which includes all $`\alpha _iP`$($`iI`$). We assume that a $``$-valued symmetric bilinear form on $`P`$ $`(,):P\times P`$ is defined and that the simple coroots $`h_iP^{}`$($`iI`$) are given as data. The natural pairing $`,:P^{}\times P`$ between $`P`$ and $`P^{}`$ is assumed to satisfy
$$h_i,\alpha _j=\{\begin{array}{cc}2\hfill & \text{if }i=j\text{ and }iI_{\mathrm{even}}\text{,}\hfill \\ 0\text{or}2\hfill & \text{if }i=j\text{ and }iI_{\mathrm{odd}}\text{,}\hfill \\ 0\hfill & \text{if }ij\text{.}\hfill \end{array}$$
We denote by $`\mathrm{\Pi }^{}=\{h_i|iI\}`$ the set of all coroots. Furthermore, for each $`iI`$ we assume that there exists a nonzero integer $`\mathrm{}_i`$ such that $`\mathrm{}_ih_i,\lambda =(\alpha _i,\lambda )`$ for every $`\lambda P`$. Then we immediately have the Cartan matrix $`A=[h_i,\alpha _j]_{ij}`$ is symmetrizable because $`\mathrm{}_ih_i,\alpha _j=(\alpha _i,\alpha _j)=(\alpha _j,\alpha _i)=\mathrm{}_jh_j,\alpha _i`$. We mention that the symmetrized matrix is $`A^{\mathrm{sym}}=\mathrm{diag}(\mathrm{}_1,\mathrm{},\mathrm{}_r)A=[(\alpha _i,\alpha _j)]_{ij}`$. Let $`𝔥=P^{}_{}`$. Then $`\mathrm{\Phi }=(𝔥,\mathrm{\Pi }^{},\mathrm{\Pi })`$ is said to be a fundamental root data associated to $`A`$. Let $`𝔤=𝔤(\mathrm{\Phi })`$ be the contragredient Lie superalgebra obtained from $`\mathrm{\Phi }`$ and $`p`$. According to , we define the quantized enveloping algebra $`U_q(𝔤)`$ to be the unital associative algebra over $`(q)`$ with generators $`q^h(hP^{}),e_i,f_i(iI)`$, which satisfy the following defining relations (compare and ):
1. $`q^h=1`$ for $`h=0`$,
2. $`q^{h_1}q^{h_2}=q^{h_1+h_2}`$ for $`h_1,h_2P^{}`$,
3. $`q^he_i=q^{h,\alpha _j}e_iq^h`$ for $`hP^{}`$ and $`iI`$,
4. $`q^hf_i=q^{h,\alpha _j}f_iq^h`$ for $`hP^{}`$ and $`iI`$,
5. $`[e_i,f_j]=\delta _{ij}{\displaystyle \frac{q^{\mathrm{}_ih_i}q^{\mathrm{}_ih_i}}{q^\mathrm{}_iq^\mathrm{}_i}}`$ for $`i,jI`$,
where $`[e_i,f_j]`$ means the supercommutator
$$[e_i,f_j]=e_if_j(1)^{p(i)p(j)}f_je_i.$$
We assume further conditions (bitransitivity condition, see p.19):
1. If $`a_{iI}U_q(𝔫_+)e_iU_q(𝔫_+)`$ satisfies $`f_iaU_q(𝔫_+)f_i`$ for all $`iI`$, then $`a=0`$,
2. If $`a_{iI}U_q(𝔫_{})f_iU_q(𝔫_{})`$ satisfies $`e_iaU_q(𝔫_{})e_i`$ for all $`iI`$, then $`a=0`$,
where $`U_q(𝔫_+)`$ (resp.$`U_q(𝔫_{})`$) is the subalgebra of $`U_q(𝔤)`$ generated by $`\{e_i|iI\}`$ (resp. $`\{f_i|iI\}`$). $`U_q(𝔤)`$ is a Hopf superalgebra whose comultiplication $`\mathrm{}`$, counit $`\epsilon `$, antipode $`S`$ are as follows.
$$\begin{array}{cc}& \mathrm{}(q^h)=q^hq^h\text{for }hP^{},\hfill \\ & \mathrm{}(e_i)=e_iq^{\mathrm{}_ih_i}+1e_i\text{for }iI,\hfill \\ & \mathrm{}(f_i)=f_i1+q^{\mathrm{}_ih_i}f_i\text{for }iI,\hfill \\ & \epsilon (q^h)=1\text{for }hP^{},\epsilon (e_i)=\epsilon (f_i)=0\text{for }iI,\hfill \\ & S(q^{\pm h})=q^h\text{for }hP^{},\hfill \\ & S(e_i)=e_iq^{\mathrm{}_ih_i},S(f_i)=q^{\mathrm{}_ih_i}f_i\text{for }iI.\hfill \end{array}$$
This is not a Hopf algebra. In order to give a Hopf algebra structure to $`U_q(𝔤)`$, we define an involutive operator $`\sigma `$ on $`U_q(𝔤)`$ by $`\sigma (q^h)=q^h`$ for all $`hP^{}`$ and $`\sigma (e_i)=(1)^{p(i)}e_i`$, $`\sigma (f_i)=(1)^{p(i)}f_i`$ for all $`iI`$. Let $`U_q^\sigma (𝔤)=U_q(𝔤)U_q(𝔤)\sigma `$. Then $`U_q^\sigma (𝔤)`$ is the algebra with the additional multiplication law given by $`\sigma ^2=1`$ and $`\sigma ^1x\sigma =\sigma (x)`$ for any $`xU_q(𝔤)`$. $`U_q(𝔤)`$ is a Hopf algebra whose comultiplication $`\mathrm{}_\sigma `$, counit $`\epsilon _\sigma `$, antipode $`S_\sigma `$ are as follows.
$$\begin{array}{cc}& \mathrm{}_\sigma (\sigma )=\sigma \sigma ,\hfill \\ & \mathrm{}_\sigma (q^h)=q^hq^h\text{for }hP^{},\hfill \\ & \mathrm{}_\sigma (e_i)=e_iq^{\mathrm{}_ih_i}+\sigma ^{p(i)}e_i\text{for }iI,\hfill \\ & \mathrm{}_\sigma (f_i)=f_i1+\sigma ^{p(i)}q^{\mathrm{}_ih_i}f_i\text{for }iI,\hfill \\ & \epsilon _\sigma (\sigma )=\epsilon _\sigma (q^h)=1\text{for }hP^{},\epsilon _\sigma (e_i)=\epsilon _\sigma (f_i)=0\text{for }iI,\hfill \\ & S_\sigma (\sigma )=\sigma ,S_\sigma (q^{\pm h})=q^h\text{for }hP^{},\hfill \\ & S_\sigma (e_i)=\sigma ^{p(i)}e_iq^{\mathrm{}_ih_i},S_\sigma (f_i)=\sigma ^{p(i)}q^{\mathrm{}_ih_i}f_i\text{for }iI.\hfill \end{array}$$
The quantized enveloping algebra $`U_q\left(𝔤𝔩(m,n)\right)`$ is obtained from the fundamental root data as follows.
* $`I=I_{\mathrm{even}}I_{\mathrm{odd}}`$ is defined by $`I_{\mathrm{even}}=\{1,2,\mathrm{},m1,m+1,\mathrm{},m+n1\}`$ and $`I_{\mathrm{odd}}=\{m\}`$,
* $`P=_{bB}ϵ_b`$, where $`B=B_+B_{}`$ with $`B_+=\{1,\mathrm{},m\}`$ and $`B_{}=\{m+1,\mathrm{},m+n\}`$,
* $`(,):P\times P`$ is the symmetric bilinear form on $`P`$ defined by
$$(ϵ_a,ϵ_a^{})=\{\begin{array}{cc}1\hfill & \text{if }a=a^{}B_+\text{,}\hfill \\ 1\hfill & \text{if }a=a^{}B_{}\text{,}\hfill \\ 0\hfill & \text{otherwise,}\hfill \end{array}$$
* $`\mathrm{\Pi }=\{\alpha _i|iI\}`$ is defined by $`\alpha _i=ϵ_iϵ_{i+1}`$,
* $`\mathrm{\Pi }^{}=\{h_i|iI\}`$ is uniquely determined by the formula $`\mathrm{}_ih_i,\lambda =(\alpha _i,\lambda )`$ for any $`\lambda P`$,
where
$$\mathrm{}_i=\{\begin{array}{cc}1\hfill & \text{if }i=1,\mathrm{},m\text{,}\hfill \\ 1\hfill & \text{if }i=m+1,\mathrm{},m+n1\text{.}\hfill \end{array}$$
Let $`V`$ be as in section2, and suppose $`R=(q)`$. The vector representation ($`\rho ,V`$) of $`U_q^\sigma \left(𝔤𝔩(m,n)\right)`$ on $`_2`$-graded vector space $`V=V_{\overline{0}}V_{\overline{1}}`$ (recall that $`V_{\overline{0}}=_{i=1}^mRv_i,V_{\overline{1}}=_{i=m+1}^{m+n}Rv_i`$) is defined by
$$\begin{array}{cc}& \rho (\sigma )v_j=(1)^{|v_j|}v_j\text{for }j=1,\mathrm{},m+n,\hfill \\ & \rho (q^h)v_j=q^{ϵ_j(h)}v_j\text{for }hP^{},j=1,\mathrm{},m+n,\hfill \\ & \rho (e_j)v_{j+1}=v_j\text{for }j=1,\mathrm{},m+n1,\hfill \\ & \rho (f_j)v_j=v_{j+1}\text{for }j=1,\mathrm{},m+n1,\hfill \\ & \text{otherwise }0\text{.}\hfill \end{array}$$
(3.1)
This representation can be extended to the representation on the tensor space $`V^r`$. Let $`\rho _r`$ be the map from $`U_q^\sigma \left(𝔤𝔩(m,n)\right)`$ to $`\mathrm{End}_{(q)}(V^r)`$ defined by
$$\begin{array}{cc}& \rho _r(\sigma )=\rho (\sigma )^r,\hfill \\ & \rho _r(q^h)=\rho (q^h)^r\text{for }hP^{},\hfill \\ & \rho _r(e_i)=\underset{k=1}{\overset{N}{}}\rho (\sigma ^{p(i)})^{k1}\rho (e_i)\rho (q^{\mathrm{}_ih_i})^{rk}\text{for }iI,\hfill \\ & \rho _r(f_i)=\underset{k=1}{\overset{r}{}}\rho (\sigma ^{p(i)}q^{\mathrm{}_ih_i})^{k1}\rho (f_i)\mathrm{Id}^{rk}\text{for }iI.\hfill \end{array}$$
(3.2)
###### Proposition 3.1 ( Proposition3.1).
$`\rho _r`$ gives a completely reducible representation of $`U_q^\sigma \left(𝔤𝔩(m,n)\right)`$ on $`V^r`$ for $`r1`$.
Making use of the comultiplication of Hopf algebra, we obtain $`\rho _r`$. Let $`\mathrm{}^{(1)}=\mathrm{}_\sigma `$ at first and set $`\mathrm{}^{(k)}=(\mathrm{}_\sigma \mathrm{Id}^{k1})\mathrm{}^{(k1)}`$ inductively. Then from the definition of $`\mathrm{}_\sigma `$, we have $`\rho _r(x)=\rho ^r\mathrm{}^{(r1)}(x)`$ for $`xU_q^\sigma \left(𝔤𝔩(m,n)\right)`$ immediately.
## 4 Schur-Weyl reciprocity for $`U_q^\sigma \left(𝔤𝔩(m,n)\right)`$ and $`_{(q),r}(q)`$
In the preceding sections, we obtained the left action of $`U_q^\sigma \left(𝔤𝔩(m,n)\right)`$ and the right one of $`_{(q),r}(q)`$ on the tensor space. We notice that the element $`q`$ of $`(q)`$ is an indeterminate. In this section we derive the commutativity between these two actions. Furthermore, we establish Schur-Weyl reciprocity for these two algebras. We consider the relation between the vector representation $`\rho _2`$ and the operator $`T`$, both act on the tensor space $`V^2`$, at first.
###### Proposition 4.1.
$`T`$ commutes with the action of $`U_q^\sigma \left(𝔤𝔩(m,n)\right)`$ on $`V^2`$ which is given by the representation $`\rho _2`$.
###### Proof.
The action $`\rho _2(g)=(\rho \rho )\mathrm{}_\sigma (g)`$ of $`U_q^\sigma \left(𝔤𝔩(m,n)\right)`$ on $`V^2`$ is defined in several cases depending upon the generator of $`U_q^\sigma \left(𝔤𝔩(m,n)\right)`$ and the basis vector $`v_iv_j`$ of $`V^2`$.
Case1 : $`g=\sigma `$
One can immediately check the commutativity between $`\rho _2`$ and $`T`$.
Case2 : $`g=q^h`$($`hP^{}`$)
It is clear in this case because of the equation,
$$\rho _2(q^h)v_iv_j=q^{ϵ_i(h)}q^{ϵ_j(h)}v_iv_j=q^{(ϵ_i+ϵ_j)(h)}v_iv_j.$$
Case3 : $`g=e_k`$ ($`k=1,2,\mathrm{},m+n`$)
We obtain the case-by-case definition of $`\rho _2`$ from (3.2) as follows.
$$\rho _2(e_k)v_iv_j=\{\begin{array}{cc}0\hfill & \text{if }i,jk+1,\hfill \\ v_{i1}v_j\hfill & \text{if }i=k+1,jk,k+1,\hfill \\ q^{(1)^{|v_j|}}v_{i1}v_j\hfill & \text{if }i=k+1,j=k,\hfill \\ q^{(1)^{|v_j|}}v_{i1}v_j\hfill & \\ +(1)^{|v_i|\delta _{km}}v_iv_{j1}\hfill & \text{if }i=k+1,j=k+1,\hfill \\ (1)^{|v_i|\delta _{km}}v_iv_{j1}\hfill & \text{if }ik+1,j=k+1.\hfill \end{array}$$
(4.1)
The operator $`T`$ has already defined in (2.1). One can check the commutativity between $`T`$ and $`\rho _2`$ by a direct computation in each case.
Case3–1 : $`i,jk+1`$
In this case, we imediately see $`\left(\rho _2(e_k)v_iv_j\right)T=\rho _2(e_k)(v_iv_jT)=0`$.
Case3–2 : $`i=k+1,j>k+1`$
$$\begin{array}{cc}\hfill \left\{\rho _2(e_k)v_iv_j\right\}T& =v_{i1}v_jT\hfill \\ & =(1)^{|v_{i1}||v_j|}v_jv_{i1}+(qq^1)v_{i1}v_j\hfill \\ \hfill \rho _2(e_k)(v_iv_jT)& =\rho _2(e_k)\left\{(1)^{|v_i||v_j|}v_jv_i+(qq^1)v_iv_j\right\}\hfill \\ & =(1)^{(|v_i|+\delta _{km})|v_j|}v_jv_{i1}+(qq^1)v_{i1}v_j\hfill \end{array}$$
If $`k<m`$, then $`|v_{i1}|=|v_i|=\delta _{km}=0`$. If $`k=m`$, then $`|v_{i1}|=0,|v_i|=1,\delta _{km}=1`$. If $`k>m`$, then $`|v_{i1}|=|v_i|=1,\delta _{km}=0`$. In each case, $`\left\{\rho _2(e_k)v_iv_j\right\}T=\rho _2(e_k)(v_iv_jT)`$ holds.
Case3–3 : $`i=k+1,j<k`$
$$\begin{array}{cc}\hfill \left\{\rho _2(e_k)v_iv_j\right\}T& =v_{i1}v_jT\hfill \\ & =(1)^{|v_{i1}||v_j|}v_jv_{i1}\hfill \\ \hfill \rho _2(e_k)(v_iv_jT)& =\rho _2(e_k)\left\{(1)^{|v_i||v_j|}v_jv_i\right\}\hfill \\ & =(1)^{(|v_i|+\delta _{km})|v_j|}v_jv_{i1}\hfill \end{array}$$
As in the case3–2, $`\left\{\rho _2(e_k)v_iv_j\right\}T=\rho _2(e_k)(v_iv_jT)`$ holds.
Case3–4 : $`i=k+1,j=k`$
$$\begin{array}{cc}\hfill \left\{\rho _2(e_k)v_iv_j\right\}T& =q^{(1)^{|v_j|}}v_{i1}v_jT\hfill \\ & =q^{(1)^{|v_j|}}2^1\left\{(1)^{|v_j|}(q+q^1)+qq^1\right\}v_jv_j\hfill \\ \hfill \rho _2(e_k)(v_iv_jT)& =\rho _2(e_k)\left\{(1)^{|v_i||v_j|}v_jv_i\right\}\hfill \\ & =(1)^{(|v_{j+1}|+\delta _{km})|v_j|}v_jv_j\hfill \end{array}$$
If $`k<m`$, then $`|v_j|=\delta _{km}=0`$. If $`k=m`$, then $`|v_j|=0,\delta _{km}=1`$. If $`k>m`$, then $`|v_j|=|v_{j+1}|=1,\delta _{km}=0`$. In each case, $`\left\{\rho _2(e_k)v_iv_j\right\}T=\rho _2(e_k)(v_iv_jT)`$ holds.
Case3–5 : $`i=k+1,j=k+1`$
$$\begin{array}{cc}\hfill \left\{\rho _2(e_k)v_iv_j\right\}T& =\left\{q^{(1)^{|v_j|}}v_{i1}v_j+(1)^{|v_i|\delta _{km}}v_iv_{j1}\right\}T\hfill \\ & =q^{(1)^{|v_i|}}(1)^{|v_{i1}||v_i|}v_iv_{i1}+\left\{q^{(1)^{|v_i|}}(qq^1)+(1)^{(|v_{i1}|+\delta _{km})|v_i|}\right\}v_{i1}v_i\hfill \\ \hfill \rho _2(e_k)(v_iv_jT)& =\rho _2(e_k)\left\{\frac{(1)^{|v_i|}(q+q^1)+qq^1}{2}v_iv_i\right\}\hfill \\ & =\left\{\frac{(1)^{|v_i|}(q+q^1)+qq^1}{2}\right\}\left\{q^{(1)^{|v_i|}}v_{i1}v_i+(1)^{|v_i|\delta _{km}}v_iv_{i1}\right\}\hfill \end{array}$$
If $`k<m`$, then $`|v_{i1}|=|v_i|=\delta _{km}=0`$. If $`k=m`$, then $`|v_{i1}|=0,|v_i|=\delta _{km}=1`$. If $`k>m`$, then $`|v_{i1}|=|v_i|=1,\delta _{km}=0`$. In each case, $`\left\{\rho _2(e_k)v_iv_j\right\}T=\rho _2(e_k)(v_iv_jT)`$ holds.
Case3–6 : $`i<k,j=k+1`$
$$\begin{array}{cc}\hfill \left\{\rho _2(e_k)v_iv_j\right\}T& =(1)^{|v_i|\delta _{km}}v_iv_{j1}T\hfill \\ & =(1)^{(|v_{j1}|+\delta _{km})|v_i|}v_{j1}v_i+(1)^{|v_i|\delta _{km}}(qq^1)v_iv_{j1}\hfill \\ \hfill \rho _2(e_k)(v_iv_jT)& =\rho _2(e_k)\left\{(1)^{|v_i||v_j|}v_jv_i+(qq^1)v_iv_j\right\}\hfill \\ & =(1)^{|v_i||v_j|}v_{j1}v_i+(qq^1)(1)^{|v_i|\delta _{km}}v_iv_{j1}\hfill \end{array}$$
If $`k<m`$, then $`|v_i|=|v_{j1}|=|v_j|=\delta _{km}=0`$. If $`k=m`$, then $`|v_i|=|v_{j1}|=0,|v_j|=\delta _{km}=1`$. If $`k>m`$, then $`|v_{j1}|=|v_j|=1,\delta _{km}=0`$. In each case, $`\left\{\rho _2(e_k)v_iv_j\right\}T=\rho _2(e_k)(v_iv_jT)`$ holds.
Case3-7 : $`i=k,j=k+1`$
$$\begin{array}{cc}\hfill \left\{\rho _2(e_k)v_iv_j\right\}T& =(1)^{|v_i|\delta _{km}}v_iv_{j1}T\hfill \\ & =(1)^{|v_i|\delta _{km}}\frac{(1)^{|v_i|}(q+q^1)+qq^1}{2}v_iv_i\hfill \\ \hfill \rho _2(e_k)(v_iv_jT)& =\rho _2(e_k)\left\{(1)^{|v_i||v_j|}v_jv_i+(qq^1)v_iv_j\right\}\hfill \\ & =\{(1)^{|v_i||v_j|}q^{(1)|v_i|}+(qq^1)(1)^{|v_i|\delta _{km}}v_iv_i\hfill \end{array}$$
If $`k<m`$, then $`|v_i|=|v_j|=\delta _{km}=0`$. If $`k=m`$, then $`|v_i|=0,|v_j|=\delta _{km}=1`$. If $`k>m`$, then $`|v_i|=|v_j|=1,\delta _{km}=0`$. In each case, $`\left\{\rho _2(e_k)v_iv_j\right\}T=\rho _2(e_k)(v_iv_jT)`$ holds.
Case3-8 : $`i>k+1,j=k+1`$
$$\begin{array}{cc}\hfill \left\{\rho _2(e_k)v_iv_j\right\}T& =(1)^{|v_i|\delta _{km}}v_iv_{j1}T\hfill \\ & =(1)^{(|v_{j1}|+\delta _{km})|v_i|}v_{j1}v_i\hfill \\ \hfill \rho _2(e_k)(v_iv_jT)& =\rho _2(e_k)(1)^{|v_i||v_j|}v_jv_i\hfill \\ & =(1)^{|v_i||v_j|}v_{j1}v_i\hfill \end{array}$$
If $`k<m`$, then $`|v_{j1}|=|v_j|=\delta _{km}=0`$. If $`k=m`$, then $`|v_{j1}|=0,|v_j|=\delta _{km}=1`$. If $`k>m`$, then $`|v_{j1}|=|v_j|=1,\delta _{km}=0`$. In each case, $`\left\{\rho _2(e_k)v_iv_j\right\}T=\rho _2(e_k)(v_iv_jT)`$ holds.
Case3-1 to Case3-8 exhaust the possible cases in (4.1), thus we have checked the commutativity for case3.
Case4 : $`g=f_k`$ ($`k=1,2,\mathrm{},m+n`$)
In a similar manner to case3, we obtain
$$\rho _2(f_k)v_iv_j=\{\begin{array}{cc}0\hfill & \text{if }i,jk,\hfill \\ (1)^{|v_i|\delta _{km}}v_iv_{j+1}\hfill & \text{if }ik,k+1,j=k,\hfill \\ v_{i+1}v_j\hfill & \\ +(1)^{|v_i|\delta _{km}}q^{(1)^{|v_i|}}v_iv_{j+1}\hfill & \text{if }i=k,j=k,\hfill \\ (1)^{|v_i|\delta _{km}}q^{(1)^{|v_i|}}v_iv_{j+1}\hfill & \text{if }i=k+1,j=k,\hfill \\ v_{i+1}v_j\hfill & \text{if }i=k,jk,\hfill \end{array}$$
and one can check for this case, so we omit the detail.
Finally, these exhaust the entirely possible cases, thus we have completed the proof. ∎
###### Proposition 4.2.
For every $`g_{(q),r}(q)`$ and $`xU_q^\sigma \left(𝔤𝔩(m,n)\right)`$, we have $`\pi _r(g)\rho _r(x)=\rho _r(x)\pi _r(g)`$.
###### Proof.
When $`r=2`$, we have already shown in proposition 4.1. It is enough to prove for $`g\{T_1,\mathrm{},T_{r1}\}`$. We deduce from cocomutativity of $`\mathrm{}_\sigma `$ that
$$\begin{array}{cc}\hfill \mathrm{}^{(r)}& =(\mathrm{}_\sigma \mathrm{id}^{r1})\mathrm{}^{(r1)}\hfill \\ & =(\mathrm{id}\mathrm{}_\sigma \mathrm{id}^{r2})\mathrm{}^{(r1)}\hfill \\ & =(\mathrm{id}^2\mathrm{}_\sigma \mathrm{id}^{r3})\mathrm{}^{(r1)}\hfill \\ & =\mathrm{}\hfill \\ & =(\mathrm{id}^{r1}\mathrm{}_\sigma )\mathrm{}^{(r1)}\hfill \end{array}$$
where $`\mathrm{id}`$ is the identity operator on $`U_q^\sigma \left(𝔤𝔩(m,n)\right)`$. For any $`r`$ with $`r>2`$, we may write $`\mathrm{}^{(r2)}(x)=x_1\mathrm{}x_{r1}`$ for some $`x_1,\mathrm{},x_{r1}U_q^\sigma \left(𝔤𝔩(m,n)\right)`$. Then applying the case $`r=2`$, we have the following.
$$\begin{array}{cc}\hfill \pi _r(T_i)\rho _r(x)& =\left\{\mathrm{Id}^{i1}\pi _2(T_i)\mathrm{Id}^{ri1}\right\}\left\{\rho ^r\mathrm{}^{(r1)}(x)\right\}\hfill \\ & =\left\{\mathrm{Id}^{i1}\pi _2(T_i)\mathrm{Id}^{ri1}\right\}\left\{\rho ^r(\mathrm{id}^{i1}\mathrm{}_\sigma \mathrm{id}^{ri1})\mathrm{}^{(r2)}(x)\right\}\hfill \\ & =\left\{_{k=1}^{i1}\rho (x_k)\right\}\pi _2(T_i)\mathrm{}_\sigma (x_i)\left\{_{l=i+1}^{r1}\rho (x_l)\right\}\hfill \\ & =\left\{_{k=1}^{i1}\rho (x_k)\right\}\mathrm{}_\sigma (x_i)\pi _2(T_i)\left\{_{l=i+1}^{r1}\rho (x_l)\right\}\hfill \\ & =\left\{\rho ^r(\mathrm{id}^{i1}\mathrm{}_\sigma \mathrm{id}^{ri1})\mathrm{}^{(r2)}(x)\right\}\left\{\mathrm{Id}^{i1}\pi _2(T_i)\mathrm{Id}^{ri1}\right\}\hfill \\ & =\rho _r(x)\pi _r(T_i)\hfill \end{array}$$
We may define $`𝒢_q`$ to be the subalgebra of $`\mathrm{End}_{R_0}\left((R_{0}^{}{}_{}{}^{m+n})^r\right)\mathrm{Mat}((m+n)^r,R_0)`$ generated by the set $`\{\rho _r(\sigma ),\rho _r(q^h),\rho _r(e_i)\rho _r(f_i)|hP^{},iI\}`$ because all the matrix elements of those generators are in $`R_0`$ from (3.1) and (3.2). For the same reason we may also define $`𝒮_q`$ the one generated by $`\{\pi _r(T_j)|j=1,\mathrm{},r1\}`$. Let us define two subalgebras of $`\mathrm{Mat}((m+n)^r,R_0)`$ as follows.
$$\begin{array}{cc}\hfill \stackrel{~}{𝒮}_q& =\{X\mathrm{Mat}((m+n)^r,R_0)|XY=YX\text{ for all }Y𝒮_q\}\hfill \\ \hfill \stackrel{~}{𝒢}_q& =\{X\mathrm{Mat}((m+n)^r,R_0)|XY=YX\text{ for all }Y𝒢_q\}\hfill \end{array}$$
Let $`R_0=[q,q^1]`$ be the Laurent polynomial ring. Let us define the specialization to a nonzero complex number $`t`$ to be a ring homomorphism $`\phi _t:R_0`$ with the condition $`\phi _t(q)=t`$. $``$ becomes $`(,R_0)`$-bimodule, with $`R_0`$ acting from the right via $`\phi _t`$. Applying the spacialization $`\phi _t`$, we obtain the specialized algebras $`𝒢_t=_{R_0}𝒢_q`$ and $`𝒮_t=_{R_0}𝒮_q`$ which are subalgebras of $`\mathrm{Mat}((m+n)^r,)`$. We also have $`\stackrel{~}{𝒢}_t=_{R_0}\stackrel{~}{𝒢}_q`$ and $`\stackrel{~}{𝒮}_t=_{R_0}\stackrel{~}{𝒮}_q`$ immediately.
###### Proposition 4.3.
$`𝒢_q=\stackrel{~}{𝒮}_q`$ and $`𝒮_q=\stackrel{~}{𝒢}_q`$ hold.
###### Proof.
Since $`R_0`$ is a principal ideal domain, the submodules $`𝒢_q`$ and $`𝒮_q`$ of the free $`R_0`$-module $`\mathrm{Mat}((m+n)^r,R_0)`$ are also free. Let $`X_i(q)\mathrm{Mat}((m+n)^r,R_0)`$($`i=1,\mathrm{},N`$) be a basis of $`𝒢_q`$ and $`x_i^{k,l}(q)R_0`$ the $`(k,l)`$-entry of $`X_i(q)`$. Then we immediately have that the specialized elements $`X_i(t)`$($`i=1,\mathrm{},N`$) generate $`𝒢_t`$ and $`dim_{}𝒢_t\mathrm{rank}_{R_0}𝒢_q`$. Because $`X_i(q)`$ are linearly independent, $`_{i=1}^N\alpha _i(q)x_i^{k,l}(q)=0`$ for $`\alpha _1(q),\mathrm{},\alpha _N(q)R_0`$ and for all $`k,l`$ imply $`\alpha _1(q)=\mathrm{}=\alpha _N(q)=0`$. We consider the system of linear equations as follows.
$$\left[\begin{array}{cccc}x_1^{1,1}(q)& x_2^{1,1}(q)& \mathrm{}& x_N^{1,1}(q)\\ x_1^{1,2}(q)& x_2^{1,2}(q)& \mathrm{}& x_N^{1,2}(q)\\ & \mathrm{}\mathrm{}\mathrm{}\\ x_1^{(m+n)^r,(m+n)^r1}(q)& x_2^{(m+n)^r,(m+n)^r1}(q)& \mathrm{}& x_N^{(m+n)^r,(m+n)^r1}(q)\\ x_1^{(m+n)^r,(m+n)^r}(q)& x_2^{(m+n)^r,(m+n)^r}(q)& \mathrm{}& x_N^{(m+n)^r,(m+n)^r}(q)\end{array}\right]\left[\begin{array}{c}\alpha _1(q)\\ \alpha _2(q)\\ \mathrm{}\\ \alpha _N(q)\end{array}\right]=\left[\begin{array}{c}0\\ 0\\ \mathrm{}\\ 0\\ 0\end{array}\right]$$
This has only the trivial solution. But applying the specialization $`\phi _t`$ to the above system, we possibly obtain a non-trivial solution. If there exists a non-trivial solution, then $`t`$ must be a zero of some Laurent polynomial whose coefficients are in $``$. Therefore if $`t`$ is a transcendental number, the above system has only the trivial solution and hence $`dim_{}𝒢_t=\mathrm{rank}_{R_0}𝒢_q`$. In the same manner, we also have that if $`t`$ is a transcendental number, then $`dim_{}𝒮_t=\mathrm{rank}_{R_0}𝒮_q`$.
We shall show that $`\mathrm{rank}_{R_0}\stackrel{~}{𝒮}_q=dim_{}\stackrel{~}{𝒮}_t`$. One can readily see that $`\mathrm{rank}_{R_0}\stackrel{~}{𝒮}_qdim_{}\stackrel{~}{𝒮}_t`$ where $`\stackrel{~}{𝒮}_t`$ is the specialized algebra of $`\stackrel{~}{𝒮}_q`$. Assume that $`X(q)=\left(x^{k,l}(q)\right)\mathrm{Mat}((m+n)^r,R_0)\stackrel{~}{𝒮}_q`$. Then $`X(q)`$ commutes with $`\pi _r(T_i)`$ for all $`i`$, hence the matrix elements $`x^{k,l}(q)`$ ($`k,l=1,\mathrm{},(m+n)^r`$) satisfy linear equations of coefficients in $`R_0`$. In the same manner as $`𝒢_q`$, one can find that the commutativity condition turns out to be the condition of solubilities of certain Laurent polynomials of coefficients in $``$. Thus if $`t`$ is a transcendental number, only the trivial equation exists, so $`\mathrm{rank}_{R_0}\stackrel{~}{𝒮}_q=dim_{}\stackrel{~}{𝒮}_t`$ holds. In a similar manner, we also have that if $`t`$ is a transcendental number, then $`\mathrm{rank}_{R_0}\stackrel{~}{𝒢}_q=dim_{}\stackrel{~}{𝒢}_t`$. Now we fix a transcendental number $`t`$. In the preceding work such like , , it has already been shown that $`\stackrel{~}{𝒮}_1=𝒢_1`$. From this fact and the inequality
$$dim_{}𝒢_1dim_{}𝒢_tdim_{}\stackrel{~}{𝒮}_tdim_{}\stackrel{~}{𝒮}_1,$$
we deduce that $`dim_{}𝒢_t=dim_{}\stackrel{~}{𝒮}_t`$. Because $`𝒢_t\stackrel{~}{𝒮}_t`$, we obtain $`𝒢_t=\stackrel{~}{𝒮}_t`$. It is known that the specialized algebra $`_{,r}(t)=_{R_0}_{R_0,r}(q)`$ is (split)semisimple. Therefore applying the double centralizer theorem, we also obtain $`𝒮_t=\stackrel{~}{𝒢}_t`$. Using the properties,
$$\begin{array}{cc}\hfill \mathrm{rank}_{R_0}𝒢_q& =dim_{}𝒢_t,\hfill \\ \hfill \mathrm{rank}_{R_0}𝒮_q& =dim_{}𝒮_t,\hfill \\ \hfill \mathrm{rank}_{R_0}\stackrel{~}{𝒢}_q& =dim_{}\stackrel{~}{𝒢}_t,\hfill \\ \hfill \mathrm{rank}_{R_0}\stackrel{~}{𝒮}_q& =dim_{}\stackrel{~}{𝒮}_t,\hfill \end{array}$$
which are already shown in the previous discussion, we readily see that $`𝒢_q=\stackrel{~}{𝒮}_q`$ and $`𝒮_q=\stackrel{~}{𝒢}_q`$. ∎
We denote $`\pi _r\left(_{(q),r}(q)\right)`$ by $`𝒜_q`$ and $`\rho _r\left(U_q^\sigma \left(𝔤𝔩(m,n)\right)\right)`$ by $`_q`$. Then we have the following.
###### Theorem 4.4.
$`\mathrm{End}__qV^r=𝒜_q`$ and $`\mathrm{End}_{𝒜_q}V^r=_q`$ hold.
###### Proof.
Obviously $`𝒜_q𝒮_q_{R_0}(q)`$ and $`_q𝒢_q_{R_0}(q)`$ as $`(q)`$-algebras. From proposition4.3 we obtain that $`\mathrm{End}__qV^r=𝒜_q`$ and $`\mathrm{End}_{𝒜_q}V^r=_q`$, and we have completed the proof. ∎
## 5 Decomposition of the tensor space
Let $`\overline{(q)}`$ be the algebraic closure of $`(q)`$. We define $`\overline{U}_q^\sigma \left(𝔤𝔩(m,n)\right)=U_q^\sigma \left(𝔤𝔩(m,n)\right)_{(q)}\overline{(q)}`$, $`\overline{𝒜}_q=𝒜_q_{(q)}\overline{(q)}`$, $`\overline{}_q=_q_{(q)}\overline{(q)}`$. Then, $`\pi _r\left(_{\overline{(q)},r}(q)\right)=\overline{𝒜}_q`$ and $`\rho _r\left(\overline{U}_q^\sigma \left(𝔤𝔩(m,n)\right)\right)=\overline{}_q`$ as $`\overline{(q)}`$-algebras of operators on $`\overline{V}^r=(V_{(q)}\overline{(q)})^r`$. Theorem4.4 still holds when we exchange the base field from $`(q)`$ to $`\overline{(q)}`$, namely, $`\mathrm{End}_{\overline{}_q}\overline{V}^r=\overline{𝒜}_q`$ and $`\mathrm{End}_{\overline{𝒜}_q}\overline{V}^r=\overline{}_q`$.
We denote by $`\mathrm{Par}(r)`$ the set of all partitions of $`r`$. By the double centralizer theorem, there is a subset $`\mathrm{\Gamma }`$ of $`\mathrm{Par}(r)`$ such that $`\overline{𝒜}_q=_{\lambda \mathrm{\Gamma }}\overline{𝒜}_{q,\lambda }`$ where $`\overline{𝒜}_{q,\lambda }`$ is the Wedderburn component corresponding to the irreducible representation of $`_{\overline{(q)},r}(q)`$ indexed by $`\lambda `$. Moreover, we also obtain the decomposition of $`_{\overline{(q)},r}(q)\overline{U}_q^\sigma \left(𝔤𝔩(m,n)\right)`$-modules,
$$\overline{V}^r=\underset{\lambda \mathrm{\Gamma }}{}H_\lambda V_\lambda ,$$
(5.1)
where $`H_\lambda `$ is the irreducible representation of $`_{\overline{(q)},r}(q)`$ indexed by $`\lambda `$, and $`V_\lambda `$ is the one of $`\overline{U}_q^\sigma \left(𝔤𝔩(m,n)\right)`$ such that $`V_\lambda V_\mu `$ if $`\lambda \mu `$. Our subject in this chapter is to determine $`\mathrm{\Gamma }`$.
Let $`H(m,n;r)=\{\lambda =(\lambda _1,\lambda _2,\mathrm{})\mathrm{Par}(r)|\lambda _jn\text{ if }j>m\}`$. Diagrams of elements of $`H(m,n;r)`$ are exactly those contained by the $`(m,n)`$-hooks. We shall show that $`\mathrm{\Gamma }=H(m,n;r)`$.
###### Theorem 5.1.
$`\overline{𝒜}_q=_{\lambda H(m,n;r)}\overline{𝒜}_{q,\lambda }`$. Hence $`\overline{V}^r=_{\lambda H(m,n;r)}H_\lambda V_\lambda `$ holds.
###### Proof.
When $`q=1`$, then Berele and Regev have already shown that
###### Theorem 5.2 (3.20 The Hook Theorem).
Let $`F`$ be an algebraic closed field of characteristic $`0`$ and $`\rho `$ the sign permutation representation on $`V^n`$ where $`V`$ is a $`(k,l)`$-dimensional vector space over $`F`$. Then
$$\rho (F[𝔖_n])=\underset{\lambda H(k,l;n)}{}A_\lambda \underset{\lambda H(k,l;n)}{}I_\lambda $$
where each $`A_\lambda `$ is the Wedderburn component corresponding to $`\lambda `$, and $`I_\lambda `$ is a simple subalgebra of $`F[𝔖_n]`$ such that $`\rho (I_\lambda )=A_\lambda `$.
Thus $`\mathrm{\Gamma }=H(m,n;r)`$ holds for $`q=1`$. Let $`t`$ be a transcendental number. We have already shown in the proof of proposition4.3 that $`dim_{}𝒢_t=dim_{}𝒢_1`$ and $`dim_{}𝒮_t=dim_{}𝒮_1`$. Let $`𝒮_t=_{\lambda \mathrm{Par}(r)}𝒮_{t,\lambda }`$ be the Wedderburn decomposition. Then, by Theorem5.2, we have $`𝒮_{1,\lambda }=0`$ if and only if $`\lambda H(m,n;r)`$. Because $`dim_{}𝒮_{t,\lambda }dim_{}𝒮_{1,\lambda }`$ for every $`\lambda \mathrm{Par}(r)`$ and $`dim_{}𝒮_t=dim_{}𝒮_1`$, we have $`dim_{}𝒮_{t,\lambda }=dim_{}𝒮_{1,\lambda }`$ for every $`\lambda \mathrm{Par}(r)`$. Thus we obtain that $`𝒮_t=_{\lambda H(m,n;r)}𝒮_{t,\lambda }`$. Since $`\overline{𝒜}_q=(𝒮_q_{R_0}(q))_{(q)}\overline{(q)}`$ and $`dim_{(q)}𝒜_q=dim_{}𝒮_t`$, we immediately get $`dim_{\overline{(q)}}\overline{𝒜}_q=dim_{}𝒮_t`$. Because $`t`$ is transcendental, $`dim_{\overline{(q)}}\overline{𝒜}_{q,\lambda }=dim_{}𝒮_{t,\lambda }`$ for every $`\lambda \mathrm{Par}(r)`$. Thus we conclude that $`\overline{𝒜}_q=_{\lambda H(m,n;r)}\overline{𝒜}_{q,\lambda }`$. The second statement is the direct consequence of the double centralizer theorem. ∎
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# Zero-bias anomaly and Kondo-assisted quasi-ballistic 2D transport
## Abstract
Nonequilibrium transport measurements in mesoscopic quasi-ballistic 2D electron systems show an enhancement in the differential conductance around the Fermi energy. At very low temperatures, such a zero-bias anomaly splits, leading to a suppression of linear transport at low energies. We also observed a scaling of the nonequilibrium characteristics at low energies which resembles electron scattering by two-state systems, addressed in the framework of two-channel Kondo model. Detailed sample-to-sample reproducibility indicates an intrinsic phenomenon in unconfined 2D systems in the low electron-density regime.
Zero-bias anomaly (ZBA) in the non-equilibrium characteristics of ballistic systems provides additional insight into the coupling mechanism of a conduction electron with its surroundings. For example, a maximum in the differential conductance ($`dI/dV_{\mathrm{SD}}`$) at the Fermi energy ($`E_\mathrm{F}`$) in case of tunnelling via magnetic impurities kondo\_tun , or through artificially confined quantum dots Kondo\_dot\_expt ; meir1 , has been explained by a Kondo-like antiferromagnetic coupling of the electron to localized impurity spin. Similar enhancement observed in clean quantum point contacts in semiconductor heterostructures has led to much controversy regarding the spin-structure and effects of lateral confinement in such systems 1D\_Kondo\_expt . On the other hand, in quasi-ballistic metallic nanobridges the ZBA shows in a local cusp-like minimum in $`dI/dV_{\mathrm{SD}}`$ at $`E_\mathrm{F}`$, which has been interpreted in terms of nonmagnetic two-channel Kondo (2CK) framework arising from the interaction of the electrons with local two-state atomic defects 2CK\_TLS\_expt . The interest in this lies in the prediction that in particle-hole symmetric case, such a model flows to a $`T=0`$ fixed point, giving rise to quantum critical behavior and non-Fermi liquid effects 2CK\_TLS\_theory ; 2CK\_matveev .
Investigation of nonequilibrium ballistic transport in spatially extended 2D systems is primarily impeded by greater scattering from background disorder. The small level spacing requires relatively larger sample dimensions, resulting in significant momentum relaxation during the course of transport. Even though recent nonequilibrium studies in mesoscopic 2D electron systems (2DES) have shown an unexpected ZBA at low background disorder self1\_ZBA , the mechanism of transport in this regime remains unclear. Here we report the experimental observation of a new and unexpected feature in the nonequilibrium quasi-ballistic 2D transport. At zero magnetic field in unconfined samples of mesoscopic dimensions, we find a strong zero-bias enhancement in $`dI/dV_{\mathrm{SD}}`$, which splits at very low $`T`$ ($`150`$ mK) by a gate voltage ($`V_\mathrm{g}`$)-dependent magnitude. This modifies the linear transport properties by introducing intermittent low-energy nonmonotonicity in the linear conductance $`G`$ ($`=dI/dV_{\mathrm{SD}}`$ at $`V_{\mathrm{SD}}=0`$, where $`V_{\mathrm{SD}}`$ is the source-drain bias), as a function of both $`T`$ and in-plane magnetic field ($`B_{||}`$). At low $`V_{\mathrm{SD}}`$, an intriguing scaling behavior of $`dI/dV_{\mathrm{SD}}`$ in $`T`$ is also observed, which indicates a 2CK-type scattering of electrons in these systems at low $`T`$.
The devices were fabricated from Si $`\delta `$-doped GaAs/AlGaAs heterostructures. In order to minimize the disorder arising from Coulomb potential of ionized dopants, we used a thick layer ($`80`$ nm) of undoped AlGaAs spacer, and adopted a slow cooling procedure self1\_ZBA to maximize the correlations in the donor layer. The as-grown mobility of the wafers range over $`\mu 13\times 10^6`$ cm<sup>2</sup>/V s, depending on the donor ($`n_\delta `$) and electron ($`n_\mathrm{s}`$) density. Over the experimental temperature range, the momentum relaxation rate was typically $`\tau ^1(0.11)\times k_\mathrm{B}T/\mathrm{}`$, indicating a quasi-ballistic nature of the transport quasi\_ballisticity . The elastic scattering length $`\lambda v_\mathrm{F}\tau 10`$ $`\mu `$m, where $`v_\mathrm{F}`$ is the Fermi velocity, provided an upper cutoff to the device dimensions, restricting the number of elastic scattering events to very few or none. Measurement of $`dI/dV_{\mathrm{SD}}`$ was carried out using standard mixed low-frequency ac/dc technique with ac excitation $`k_\mathrm{B}T`$ at each $`T`$ inside a top-loading dilution refrigerator with base electron temperature $`32`$ mK.
Fig. 1(a), (b) and (c) show the $`B_{||}=0`$ nonequilibrium characteristics of three samples with different dimensions and doping profile. We varied both the total geometrical area, as well as the aspect ratio in samples T45 ($`2\times 8\mu `$m<sup>2</sup>, $`n_\delta 2.5\times 10^{12}`$ cm<sup>-2</sup>), T46 ($`5\times 8\mu `$m<sup>2</sup>, $`n_\delta 1.9\times 10^{12}`$ cm<sup>-2</sup>) and T48 ($`5\times 5\mu `$m<sup>2</sup>, $`n_\delta 0.9\times 10^{12}`$ cm<sup>-2</sup>). The dimensions are chosen in order to ensure small single-particle level spacing with $`\mathrm{\Delta }ϵ_\mathrm{x}`$,$`\mathrm{\Delta }ϵ_\mathrm{y}h^2/8m^{}L^2k_\mathrm{B}T`$ for each $`T`$. While no consistent feature in $`G`$ as a function of $`V_\mathrm{g}`$ could be identified (not shown), the ZBA around $`V_{\mathrm{SD}}=0`$ is evident for all samples disorder . We focus on some of the common features of the ZBA shown in Fig. 1: (1) At $`|V_{\mathrm{SD}}|0.150.2`$ meV, $`dI/dV_{\mathrm{SD}}`$ shows an enhancement at all $`V_\mathrm{g}`$. This energy scale varies only weakly with $`V_\mathrm{g}`$, and also $`E_\mathrm{F}`$ at all $`n_\mathrm{s}`$. (2) Over certain ranges of $`V_\mathrm{g}`$, the ZBA splits by $`2\mathrm{\Delta }`$, leading to a double-peak structure around $`E_\mathrm{F}`$. (3) $`\mathrm{\Delta }`$ is oscillatory in $`V_\mathrm{g}`$, and when normalized for series conductance, in particular for $`G2e^2/h`$, the splitting is strongest when $`G`$ lies close to an even integral multiple of $`e^2/h`$, becoming unresolvable around odd multiples. (4) While no clear dependence of the energy scales on size/shape of the samples was observed, the amplitude of the ZBA was found to depend on $`n_\mathrm{s}`$. Note that the ZBA is strongest in T48 ($`n_\mathrm{s}=6.5\times 10^{10}`$ cm<sup>-2</sup> at $`V_\mathrm{g}=0`$), and weakest in T45 ($`n_\mathrm{s}=10.1\times 10^{10}`$ cm<sup>-2</sup> at $`V_\mathrm{g}=0`$). Such weakening is also observed in a given sample as $`V_\mathrm{g}`$ (or $`n_\mathrm{s}`$) is increased.
The single- and double-peak regions show distinct behavior in the presence of external in-plane magnetic field ($`B_{||}`$). As shown in Fig. 1(d) and (e), while the single-peak splits by the Zeeman energy $`\mathrm{\Delta }_\mathrm{Z}=2g^{}\mu _\mathrm{B}B_{||}`$ from $`B_{||}=0`$, Zeeman splitting of the double peak appears only at relatively large (sample-dependent) $`B_{||}`$-scale ($`0.52`$ Tesla)self1\_ZBA . Such a splitting of ZBA is taken as a distinctive feature of Kondo-type dynamics kondo\_tun ; Kondo\_dot\_expt . The effective $`g`$-factor $`g^{}`$ was found to be both sample- and $`n_\mathrm{s}`$-dependent $`g^{}/g_\mathrm{b}13`$, where $`|g_\mathrm{b}|=0.44`$, consistent with exchange-induced enhancement at low $`n_\mathrm{s}`$ self1\_ZBA ; g\_factor .
Kondo-type dynamics also lead to a suppression of the ZBA with increasing $`T`$. Indeed, such a suppression is observed in all our samples as shown in Fig. 2a for nine different $`V_\mathrm{g}`$’s in T46. In case of single peaks (sets 2, 3, 7 and 8), the suppression is monotonic as $`T`$ is increased from $``$ 30 mK to $``$ 800 mK. For $`\mathrm{\Delta }0`$, the $`T`$ dependence of $`dI/dV_{\mathrm{SD}}`$ is nonmonotonic in $`T`$ close to $`E_\mathrm{F}`$ (see sets 1, 4, 5, 6, and weakly in 9). Fig. 2b shows the $`T`$-dependence of $`G`$ from the traces in Fig. 2a. Qualitatively, the high-$`T`$ regime (typically $`T300`$ mK) is similar for all cases, where $`G`$ shows a “metal”-like decrease with increasing $`T`$. For $`T100150`$ mK, the increasing behavior of $`G`$ with $`T`$ appears intermittently, at values of $`V_\mathrm{g}`$ where $`dI/dV_{\mathrm{SD}}`$ shows resolvable double-peak structure (Fig. 2b and 3b-e). This reentrant nature is crucial, because even if the high-$`T`$ suppression of $`G`$ is attributed to various combinations of phonon contribution, the interaction correction, or the Altshuler-Aronov-type correction from electron scattering by Friedel oscillations int\_corr , the repeated change in the sign of $`dG/dT`$ is clearly inconsistent with standard weak-localization/interaction-based mechanisms having a single transition from localized to “metallic” state transport. However, in the Zeeman split regime at high $`B_{||}`$, $`G`$ increases with $`T`$ at all $`V_\mathrm{g}`$, indicating standard localized state transport.
Most striking aspect of the double-peak structure is the collapse of $`\mathrm{\Delta }G(V_{\mathrm{SD}},V_\mathrm{g},T)=dI/dV_{\mathrm{SD}}(V_{\mathrm{SD}},V_\mathrm{g},T)G(V_\mathrm{g},T)`$, onto a single $`V_\mathrm{g}`$-dependent trace when $`V_{\mathrm{SD}}`$ is scaled by $`T`$. In Fig. 2(c) and (d) this is illustrated for ZBA’s at two different $`G`$ (and hence $`V_\mathrm{g}`$), indicating a common underlying mechanism. The insets contain the actual data prior to substraction of $`G(V_\mathrm{g},T)`$, showing the thermal broadening of $`dI/dV_{\mathrm{SD}}`$, and the suppression of the ZBA splitting with increasing $`T`$. In Fig. 2e we show the $`T`$-dependence of individual peaks in $`dI/dV_{\mathrm{SD}}`$ in the double-peak region at the given $`V_\mathrm{g}`$. For comparison, the individual $`T`$-dependences are normalized by the $`dI/dV_{\mathrm{SD}}`$ at lowest $`T`$. Note the approximately logarithmic $`T`$-dependence of both peaks at $`T150`$ mK, which is also expected in the Kondo-framework.
The low-$`V_{\mathrm{SD}}`$ scaling of $`dI/dV_{\mathrm{SD}}`$ indicates nonequilibrium conductance of the form,
$$\frac{dI}{dV_{\mathrm{SD}}}(V_{\mathrm{SD}},V_\mathrm{g},T)=G(V_\mathrm{g},T)+AT^\alpha [V_\mathrm{g},\frac{eV_{\mathrm{SD}}}{k_\mathrm{B}T}]$$
(1)
where $`A`$ is a phenomenological constant, and $`\alpha =0`$ gives the best collapse onto the single sample-dependent function $`(V_\mathrm{g},eV_{\mathrm{SD}}/k_\mathrm{B}T)`$. The dependence of $``$ on $`V_\mathrm{g}`$ will be discussed later. The scaling form of Eq. 1 has been observed in the context of the scattering of electrons from bistable systems 2CK\_TLS\_expt ; ralph\_PRL , and addressed in a two-channel Kondo (2CK) framework 2CK\_TLS\_theory ; 2CK\_matveev . In quasi-ballistic nanostructures, 2CK dynamics may arise from two possibilities: (1) Electron scattering off systems with quasi-degenerate orbital states acting as pseudospins, while the real spins act as the channel index. The nature of this orbital degeneracy can however vary from equivalent lattice defects in metallic nanobridges ralph\_PRL , to singlet-triplet degeneracy in multilevel quantum dots kouwenhoven\_dot\_2CK ; Hofs\_Scho . (2) An underscreened high-spin ($`S`$) ground state coupled to $`M`$ ($`<2S`$) conduction channels Pust\_Glaz . When parametrized in terms of a two-impurity Kondo problem, two spins $`S_1`$ and $`S_2`$ ($`S=S_1+S_2`$) interact with conduction electrons with antiferromagnetic exchange parameters $`J_1`$ and $`J_2`$ ($`J_2J_1`$), and a direct exchange $`I`$ ($`J_1,J_2I`$vojta . A two-stage screening process can then decouple the spins from the conduction band by forming $`S_1`$-$`S_2`$ singlet if $`I`$ exceeds some critical magnitude. Since both stages can be suppressed by lifting the spin degeneracy with Zeeman energy, a distinctive feature of this mechanism is the nonmonotonicity of $`G`$ in both $`T`$ and $`B_{||}`$, with a quantitative correspondence between the respective energy scales Pust\_Glaz .
In Fig. 3, we compare the $`T`$-dependence of $`G`$ to its $`B_{||}`$-dependence in T45 (on a separate cooldown). Four representative $`V_\mathrm{g}`$, with corresponding nonequilibrium traces at $`T32`$ mK and $`B_{||}=0`$ T, are identified as $`V_{\mathrm{g1}}`$ to $`V_{\mathrm{g4}}`$ in Fig. 3a. Figs. 3b-e show the ($`B_{||}=0`$) $`T`$-dependence of $`G`$ at these $`V_\mathrm{g}`$s, while Figs. 3f-i show the $`B_{||}`$-dependence at the base $`T32`$ mK. We note that apart from the qualitative nonmonotonic behavior of $`G`$ as a function of both $`T`$ and $`B_{||}`$ at $`V_{\mathrm{g1}}`$ and $`V_{\mathrm{g3}}`$, there is also a quantitative agreement in the energy-scales over which the double-peak structures are suppressed. For example, as the half width $`\mathrm{\Delta }/2`$ reduces from $`0.03`$ meV at $`V_{\mathrm{g1}}`$ to $`0.017`$ mV at $`V_{\mathrm{g2}}`$, we find a corresponding decrease in the characteristic thermal and Zeeman energy scales (denoted by the vertical arrows) from $`0.023`$ meV and $`0.028`$ meV respectively at $`V_{\mathrm{g1}}`$, to $`0.011`$ meV and $`0.02`$ meV respectively at $`V_{\mathrm{g2}}`$. Since $`B_{||}`$ is applied in the plane of the 2DES (to an accuracy better than $`0.2^{}`$), conventional weak-localization effects are unlikely to contribute to the observed behavior GB\_WL . Coupling of $`B_{||}`$ to orbital degree of freedom through spin-orbit interaction GB\_WL\_SO or finite thickness effect GB\_WL\_thick are also excluded since neither antilocalization to weak-localization crossover, nor a stronger suppression of $`G`$ with increasing $`B_{||}`$ at lower $`n_\mathrm{s}`$ were observed.
The Kondo-type enhancement in mesoscopic 2D nonlinear transport in the ultra-clean limit indicates an unexpected scattering mechanism of the lead electrons. Moreover, intermittent splitting of the ZBA implies a two-state nature of the scatterer that becomes resolvable only at very low $`T`$ ($`150`$ mK). In view of the quasi-ballistic nature of transport, the non-equilibrium conductance can represent the tunnelling characteristics between the source and drain across the potential barrier formed by the gate. Presence of localized magnetic states within the barrier region of traditional tunnel junctions has been shown to result in Kondo-type enhancement in the tunnelling conductance kondo\_tun . In nanostructures fabricated from MBE-grown high-quality GaAs/AlGaAs heterostructure, localized acceptor sites close to the system can significantly modify the transport through capacitive coupling cobden . This could lead to a Kondo-assisted tunnelling, similar to that observed in metallic nanobridges buhrman , and also in semiconductor-metal junctions wolf . Arguments against such a scenario are, (1) reproducibility, and insensitivity of the relevant energy scales to sample-specific details, and (2) enhancement in the amplitude of ZBA at lower disorder.
Alternatively, dependence of the effect on $`n_\mathrm{s}`$, i.e. relatively strong amplitude of ZBA in samples with low-$`n_\mathrm{s}`$, and also weakening of the effect with increasing $`n_\mathrm{s}`$ within a given sample, indicates a possible role of Coulomb interaction. Typically in our systems, the low magnitude of $`n_\mathrm{s}`$ ($`0.83\times 10^{10}`$ cm<sup>-2</sup>) results in a large interaction parameter, $`r_\mathrm{s}=1/a_\mathrm{B}\sqrt{\pi n_\mathrm{s}}3.56`$, where $`a_\mathrm{B}`$ is the effective Bohr radius. This corresponds to an exchange energy that can lead to strong many-body spin fluctuations with magnetic moment $`1/2`$ many\_body\_spin . Underscreening of such spin fluctuations when Kondo-coupled to lead electrons can result in split Kondo resonance, which would be strongly $`V_\mathrm{g}`$-dependent. Exchange splitting of Kondo-resonance in presence of itinerant electron ferromagnetism has been observed experimentally pasupathy . Here, the suppression of the double peak pattern in $`B_{||}`$ provides further evidence of the dynamics to be related to real spin, rather than pseudospin resulting from quasi-degenerate orbital states. Indeed, similar ZBA observed in quasi-1D quantum point contacts have been interpreted in terms of a dynamic spin-polarization of the electrons within the 1D channel 1D\_Kondo\_expt .
Finally, we discuss the intermittent collapse of $`\mathrm{\Delta }`$ observed in Fig. 1. If the scaling relation of Eq. 1 originates from proximity of a 2CK fixed point 2CK\_TLS\_expt ; ralph\_PRL , the low-$`T`$ $`V_\mathrm{g}`$-dependence of $`\mathrm{\Delta }`$ indicates that such scaling should also be possible in terms of a $`V_\mathrm{g}`$-dependent energy-scale $`\mathrm{\Gamma }(V_\mathrm{g})`$. This implies an asymptotic form of the scaling function $`[V_{\mathrm{SD}}/\mathrm{\Gamma }(V_\mathrm{g})]`$ as $`T0`$. Theoretically, $`\mathrm{\Gamma }`$ is analogous to the energy asymmetry in case of orbital degeneracy, or the direct exchange parameter $`I`$ in the two-impurity Kondo parametrization Hofs\_Scho ; Pust\_Glaz ; vojta ; Pust\_Glaz\_2004 , which depends implicitly on $`V_\mathrm{g}`$ through the exchange interaction magnitude $`J`$ tarucha\_J . Indeed, in Fig. 4a and 4b we have shown the low-energy scaling of $`\mathrm{\Delta }G(V_{\mathrm{SD}},V_\mathrm{g})=dI/dV_{\mathrm{SD}}(V_{\mathrm{SD}},V_\mathrm{g})G(V_\mathrm{g})`$ in T46 and T48 at the base $`T32`$ mK. In the absence of an absolute scale, we have expressed $`\mathrm{\Gamma }(V_\mathrm{g})`$ with respect to its minimum value ($`\mathrm{\Gamma }(V_\mathrm{g})^{\mathrm{min}}`$) observed around $`G3.5\times e^2/h`$ in both samples (Fig. 4c and 4d). The scalability of $`\mathrm{\Delta }G`$ close to the collapse of $`\mathrm{\Delta }`$ strongly indicates the possibility of a quantum critical dynamics. However, a satisfactory explanation of its reentrant nature and the fundamental mechanism of the 2CK-type scattering forms the basis of ongoing investigations.
A.G. acknowledges fruitful discussions with G. Gumbs and V. I. Fal’ko. This work was supported by an EPSRC funded project. C.S. acknowledges financial support from Gottlieb Daimler- and Karl Benz-Foundation.
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# Phase Space Tomography of Classical and Nonclassical Vibrational States of Atoms in an Optical Lattice
## I Introduction
Much of the excitement and promise of new fields such as quantum information processing would not be possible without the development of sophisticated techniques for manipulating and measuring quantum systems. In systems such as ion traps, quantum dots, Bose-Einstein condensates, and entangled photons, it is the fact that quantum states can be accurately prepared and observed which has enabled a wide range of experimental advances, and held out hope for the creation of practical quantum technologies. In such systems, quantum states inevitably evolve into partially mixed states, which can be characterized either by a density matrix, in the case of a finite-dimensional Hilbert space, or by phase space quasi-probability distributions such as the Wigner functionwigner or the Husimi distributionhusimi . Techniques for extracting these functions, generally referred to as “quantum tomography,” have long been a topic of active research, in the context of the electromagnetic fieldraymer ; schiller , Rydberg stateswalmsley , neutral atomsjessen1 ; mlynek1 , dissociating moleculesmolmer , entangled photonswhite ; kwiat ; mitchell , and ion trapscirac1 ; wineland .
One system which has led to many interesting effects and proposals is the “optical lattice,” in which atoms are trapped in a periodic potential formed by a standing wave of light beamsprentiss ; phillips1 ; jessen2 . This analog condensed-matter system has been used to study Bragg scatteringphillips2 , precision measurement using atom interferometrykasevich , the fragmentation of Bose condensateskasfrag ; spekkens , squeezed states of atomic motionsqueezed , quantum feedbackraithel , the Mott insulator transitionbloch1 , and quantum logic gatesjessen3 ; cirac2 ; bloch2 , to name only a few examples. Typical probes of the system include the Bragg scattering of probe beams, which is sensitive to the localization $`X^2`$ of the atoms; the time-varying transmission of the lattice beams themselves, which is sensitive to the instantaneous force exerted on the atoms, related to the centre-of-mass position $`X`$ in each well; and atom-interferometric probes of long-range coherence. Here we present for the first time a complete characterisation of the quantum state of atoms trapped in optical-lattice wells, by extracting the Husimi and Wigner distributions for atoms in several different initial states. In particular, we demonstrate the “nonclassical” Wigner function for atoms with a population inversion, whose negative value at the origin is analogous to that for the 1-photon Fock state and for the excited state of the single ion in a trap.
## II Theory
The Husimi distribution is a well known quasi-probability distribution related to a state’s density matrix $`\rho `$ by the equation $`Q\left(\alpha \right)=\frac{1}{\pi }\alpha \left|\rho \right|\alpha `$, where $`|\alpha `$ is a coherent state, that is, a Gaussian distribution centered at $`\alpha `$, with minimum-uncertainty width which can be defined by a specified harmonic oscillator. It has the practical advantage that its value for a particular $`\alpha `$ is an observable and can be obtained directly with one measurement. Because it is always real and always positive it is sometimes referred to as a classical quasi-probability distribution.
The Wigner distribution, which has a one-to-one correspondence with the Husimi distribution, is particularly useful for identifying nonclassical states such as Fock or inverted states, for which it takes on negative value (as in fact it does for any non-Gaussian states). It can be defined by the expression $`W(x,p)=\frac{1}{\pi }_{\mathrm{}}^+\mathrm{}x+q\left|\rho \right|xq\mathrm{exp}\left[2ipq\right]𝑑q`$ leonhardt . The Wigner distribution is unique in that it allows one to determine the marginal probability distribution of either coordinate by integrating over the other: $`P\left(x\right)=W(x,p)𝑑p`$ and $`P\left(p\right)=W(x,p)𝑑x`$.
The measurement of these quasi-probability distributions becomes particularly simple in a harmonic oscillator. The Husimi distribution can be measured by evaluating the expression
$$Q(\left|\alpha \right|,\theta )=\frac{1}{\pi }\alpha \left|\rho \right|\alpha =\frac{1}{\pi }0|D^{}\left(\left|\alpha \right|\right)R\left(\theta \right)\rho R^{}\left(\theta \right)D\left(\left|\alpha \right|\right)|0,$$
()
where $`\theta =\mathrm{arg}\left\{\alpha \right\}`$. That is, instead of projecting the unknown state $`\rho `$ onto a coherent state $`\alpha `$, one can perform position displacement, $`D\left(\left|\alpha \right|\right)`$, and rotation, $`R\left(\theta \right)`$, operations on the unknown state and measure the overlap of the resulting state onto the ground state $`|0`$, which is a straightforward process in our experiment. Applying the displacement operator amounts to physically displacing the state a distance of $`x=2x_0\left|\alpha \right|`$, where $`x_0=\left(\frac{\mathrm{}}{2m\omega }\right)^{1/2}`$ is the ground state width of a particle of mass $`m`$ in a harmonic oscillator of frequency $`\omega `$. In a harmonic oscillator $`R\left(\theta \right)=\mathrm{exp}\left[ia^{}a\theta \right]`$ can be implemented by letting the state evolve under $`=\mathrm{}\omega \left(a^{}a+\frac{1}{2}\right)`$ for a time $`t=\theta /\omega `$. How these operations are performed experimentally and how the ground state population of a state is measured will be described in section III.
The measurement of the Wigner distribution is simplified by recognizing that for any symmetric non-degenerate potential and, in particular, in a harmonic oscillator, $`W(0,0)=\frac{1}{\pi }_{\mathrm{}}^+\mathrm{}q\left|\rho \right|q𝑑q=\frac{1}{\pi }_{n=0}^{\mathrm{}}\left(1\right)^nn\left|\rho \right|n`$. For values away from the origin the displacement and rotation operators can once again be utilized:
$$W(x,p)=W(2x_0\mathrm{}e\left(\alpha \right),2p_0\mathrm{}m\left(\alpha \right))=\frac{1}{\pi }\underset{n=0}{\overset{\mathrm{}}{}}\left(1\right)^np\left(n|\alpha \right),$$
()
where $`p_0=\sqrt{m\mathrm{}\omega /2}`$ and $`p\left(n|\alpha \right)=n|D^{}\left(\left|\alpha \right|\right)R\left(\theta \right)\rho R^{}\left(\theta \right)D\left(\left|\alpha \right|\right)|n`$ is the probability of finding the particle to be in state $`n`$ after the application of the rotation and displacement operators to the unknown state $`\rho `$. Therefore, the determination of the Wigner distribution is reduced to performing population measurements after applying displacement and rotation operators. This is similar to the reconstruction of the Husimi distribution except that the final population measurement requires the measurement of all states, ($`|0,|1,|2,\mathrm{}`$), not just the population of the ground state.
## III Implementation
### III.1 Optical Lattice Specifications
A 1-D optical lattice is formed by the counter-propagating components of two laser beams resulting in an intensity interference pattern of the form $`I\left(x\right)=I_0\mathrm{cos}^2\left(kx\mathrm{sin}\left(\frac{\gamma }{2}\right)\right)`$, where $`k=\frac{2\pi }{\lambda }`$ is the wave vector of the laser and $`\gamma `$ is the angle of intersection of the two beams. This standing light wave induces a light shift on the atoms resulting in a potential $`U\left(x\right)=I\left(x\right)\frac{\mathrm{}\mathrm{\Gamma }^2}{4\mathrm{\Delta }I_s}`$ where $`\mathrm{\Delta }`$ is the detuning of the laser light from atomic resonance and $`\mathrm{\Gamma }`$and $`I_s`$ are the natural line-width and saturation intensity of the atom respectively. Of importance to this experiment is the fact that the individual wells of an optical lattice can be approximated as harmonic oscillators. Therefore the theory of Husimi and Wigner distribution measurements in harmonic oscillators can be applied here. The oscillation frequency in each well $`\omega `$, and lattice depth $`U_0`$, are related by the equation $`\omega =\frac{4k_L}{\pi }\sqrt{\frac{U_0}{m}}`$. We find $`\omega `$ directly by measuring the period of Ramsey fringes created by inducing oscillations with a spatial shift of the sinusoidal potentialmyrskog .
Our experiment starts with a cloud of $`{}_{}{}^{85}Rb`$ atoms in a magneto optical trap (MOT) which are then cooled in an optical molasses to a temperature on the order of $`10\mu K`$. The optical lattice is turned on during the MOT stage as cooling in the presence of the lattice increases the lattice loading efficiency. The two laser beams intersect at an angle of $`\gamma =49.6^{}`$, resulting in a lattice vector of $`k_L=\frac{2\pi }{\lambda }\mathrm{sin}\left(\frac{\gamma }{2}\right)=\frac{\pi }{a}=3.3810^6m^1`$, where $`\lambda =780nm`$ is the wavelength of the lattice light and $`a=0.93\mu m`$ is the spatial period of the lattice. The lattice has a detuning of $`\mathrm{\Delta }2\pi 25GHz`$ from the $`F=3F^{}=4`$ $`D2`$ trapping line of $`{}_{}{}^{85}Rb`$ so as to make the scattering rate negligible and to allow the sinusoidal lattice potential to be treated as conservative. The lattice is tailored to support $`24`$ bound states, depending on the experiment, which means a depth of $`U_01040E_r`$ where $`E_r=\frac{\mathrm{}^2k_L^2}{2m}`$ is the scaled recoil energy in the lattice direction. The lattice is formed in the vertical direction so that atoms in unbound states will, in a time of $`10ms`$, leave the interaction region due to the pull of gravity. Each of the lattice beams passes through an acousto-optic modulator (AOM). By controlling the frequency and phase of the signal with which we drive each AOM independently we are able to control the position, velocity and acceleration of the lattice. Shifting the relative phase of the lattice beams by $`\varphi `$ displaces the lattice (or, in the rest frame of the lattice, displaces the atoms) by a distance $`d=a\varphi /2\pi `$. Spatial shifts, used for the displacement operators $`D\left(\left|\alpha \right|\right)`$, can be applied with a resolution on the order of $`1nm`$. Displacements as large as twice the lattice spatial period occur in a time of $`0.5\mu s`$. With typical oscillation frequencies of $`\omega /2\pi 10^4Hz`$, this time interval for a shift can be considered to be instantaneous.
### III.2 State Preparation
The preparation of the states whose quasi-probability distributions we shall measure requires two steps. The first is to filter out the ground state. By lowering the intensity of the lattice beams over a time of $`10ms`$ until only one bound state is supported, and keeping it there for $`5ms`$, all atoms in higher states become unbound and fall out of the interaction region due to gravity. Afterwards we raise the intensity back to its original level but are left with the ground state in each well, with a typical contamination of the first excited state of about $`515\%`$. A more detailed explanation of our ground state preparation and measurement process can be found in reference myrskog and a similar technique is described in reference mlynek2 . The implementation of the second part of the state preparation depends on what state we wish to prepare.
In this paper we prepare and measure three states: a ground state, a near coherent state and a state with a population inversion. For the ground state no further action is necessary as it was prepared in the filtering stage. A near coherent state, $`|\beta `$, is prepared by shifting the potential by $`\delta x=\frac{a}{6}=0.155\mu m`$ (or a $`60^{}`$ phase shift). The magnitude of $`\beta `$ is related to the displacement by $`\left|\beta \right|=\delta x\left(\frac{m\omega }{2\mathrm{}}\right)^{1/2}=0.88`$. For both the ground and coherent state a lattice depth of $`37E_r`$ was chosen, supporting $`4`$ bound states with an oscillation frequency of $`\omega =48.33kHz`$. The deviation of the near coherent state which we prepare from an actual coherent state is due to the finite depth of the lattice. With only $`4`$ bound states only the first $`4`$ terms of the coherent state are present in the lattice. This does not significantly detract from the validity of the approximation, as the amplitude of the $`n^{\text{th}}`$ term in a harmonic oscillator coherent state is $`\mathrm{exp}\left[\left|\beta \right|^2/2\right]\frac{1}{\sqrt{n!}}\left|\beta \right|^n`$, which when $`|\beta |=0.88`$ is quite small for $`n4`$. This state was allowed to rotate in the potential for a time $`t=20\mu s`$, giving it a rotation of $`\theta =0.97`$ radians.
Creation of an inverted state begins with the preparation of the ground state, this time in a lattice containing only two bound states ($`\omega =32.2kHz`$). Next, the potential is given a $`60^{}`$ phase shift ($`0.155\mu m`$), held there for $`80\mu s`$ and then shifted back to its original position. We have found that this process excites a large number of atoms into higher statesmaneshi . After waiting several milliseconds in order to let unbound atoms leave the interaction region we are left with what will be shown to be an incoherent mixture of ground and first-excited-state atoms in a ratio of roughly $`3`$ to $`7`$.
### III.3 State Measurement
Here we describe the process used to determine the Husimi distribution for the ground state and coherent state. As per equation 1 the measurement of $`Q(\left|\alpha \right|,\theta )`$ takes three steps. First, we allow the atoms to undergo free evolution for a time $`t`$ to let the state oscillate in the harmonic-like potential for a rotation of $`\theta =\omega t`$ in phase space. We use a total of $`27`$ angles (or wait times, $`t`$) separated by $`0.24`$ radians (or $`5\mu s`$), spanning a range of $`\theta ϵ[0,2\pi ]`$ ($`\delta tϵ[0,130]\mu s`$). Second, a displacement of $`x=2x_0\left|\alpha \right|`$ is applied to the lattice. This is the displacement operator, $`D\left(\left|\alpha \right|\right)`$. A total of $`19`$ different displacements are used, each separated by $`25.8nm`$ or a $`10^{}`$ phase shift of the potential for a total range of $`xϵ[0,465]nm`$ or $`[0,180]`$ degrees. Along with the phase component, the total number of measurements is then $`513`$ not including repeated measurements for statistical analysis and averaging. Next, a projection of this new state onto the ground state is performed. This first requires the same process that was used to filter the ground state during the state preparation: the lattice laser intensity is lowered until only the ground state remains. After waiting $`20ms`$ in order to let the unbound atoms become spatially separated from those still bound the atoms are illuminated with resonant light, and their fluorescence collected in a CCD camera. With this image the relative population of ground state atoms, as a fraction of total atom number, is measured.
Figure 1 is a phase space diagram of the Husimi distribution for the coherent state. As expected, it is a Gaussian displaced from the origin by $`\left|\beta \right|=0.88`$ at an angle of roughly $`55^{}`$. Since a coherent state is merely a ground state displaced from the origin in phase space both distributions should have the same width. Figure 2 shows a cross section of the Husimi distribution for the ground state along with one of the coherent state. Fits of these Gaussians give $`rms`$ widths of $`144nm`$ and $`147nm`$ and peaks of $`0.267`$ and $`0.257`$ for the ground and coherent states respectively. Husimi distributions, by definition, are normalized. By integrating the curve shown in figure 1 we find from$`Q(\left|\alpha \right|,\theta )\left|\alpha \right|d\left|\alpha \right|𝑑\theta =Q(\frac{r}{2x_{rms}},\theta )\frac{r}{4x_{rms}^2}𝑑r𝑑\theta =1`$ that $`x_{rms}=96.3nm`$. Here we find a surprising result. The ground state width of a harmonic oscillator of this frequency is $`x_0=\left(\frac{\mathrm{}}{2m\omega }\right)^{1/2}=88.0nm`$ but we find $`x_{rms}=96.3nm=1.09x_0`$. We believe that this discrepancy is due to inhomogeneities in the intensity of our lattice beams. We have separately measured the dephasing time $`T_2`$ to be on the order of 2 or 3 oscillation periods in our latticemyrskog ; maneshi , implying $`\mathrm{\Delta }\omega 0.4\omega `$. As $`x_0^2=\mathrm{}/2m\omega `$, we expect the rms width of our atom clouds to be $`x_{rms}=(\mathrm{}/2m)1/\omega =x_0^2[1+(\mathrm{\Delta }\omega /\omega )^2+(\mathrm{\Delta }\omega /\omega )^4+\mathrm{})]`$. For $`\mathrm{\Delta }\omega /\omega 0.4`$, this means $`x_{rms}1.08x_0`$, consistent with our observations.
In addition, the peak of the Husimi distribution of a true ground state is $`\frac{1}{\pi }`$. Our measured distribution has a peak value of $`0.267=0.84/\pi `$, obtained from a direct measurement of the ground-state population to be $`84\%`$. A harmonic-oscillator state with $`16\%`$ population in the first excited state, when the ground-state width is taken to be $`x_{rms}`$ as calculated above, has a Husimi distribution with a width of $`141`$ nm, consistent with our measured value.
Following equation 2 the Wigner distribution for the inverted state is measured by first applying the rotation and displacement operators $`R\left(\theta \right)`$ and $`D\left(\left|\alpha \right|\right)`$, as in the Husimi reconstruction. For the rotation operators we sample $`41`$ different rotation angles separated by $`0.16`$ radians spanning a range of $`\theta ϵ[0,2\pi ]`$ equivalent to waiting times in the range $`\delta tϵ[0,200]\mu s`$ with a resolution of $`5\mu s`$. We again use $`19`$ different spatial shifts for the displacement operators with a range of $`xϵ[0,465]nm`$ giving a total number of $`779`$ measurements for the Wigner measurement.
After the rotation and displacement operations are implemented, the state populations must be measured. Since the lattice holds two bound states any atoms excited to higher states by the rotation and displacement operators become unbound. After waiting $`5ms`$ in order to let these atoms become spatially separated we once again lower the lattice laser intensity until only one bound state remains. After waiting another $`20ms`$ in order to let these first-excited-state atoms leave the interaction region we then capture a fluorescence image. From this image we determine the ground and first-excited-state populations as well as how many atoms were in states $`n2`$. Therefore we are able to construct the Wigner distribution with the following caveat: measurements of $`W(x,p)`$ become increasingly uncertain for values for which the displacement operator $`D\left(\left|\alpha \right|\right)`$ excites more atoms into $`n2`$, since only the first two terms of equation 2 are known exactly.
Figure 3 shows a cross section of our best estimate of the Wigner distribution: $`W^{}(x,0)\frac{1}{\pi }\left(p\left(0|\alpha \right)p\left(1|\alpha \right)\right)`$. At the origin we see that the distribution is negative, a nonclassical signature which is characteristic of a population inversion. The inset shows the absolute upper and lower bounds, based on our data, of the Wigner distribution due to our lack of knowledge about the populations of states higher than $`n=1`$. The upper bound is obtained by assuming that all atoms lost during the displacement and rotation operators were in the 2nd excited state; $`\frac{1}{\pi }\left(p\left(0|\alpha \right)p\left(1|\alpha \right)+p\left(n>1|\alpha \right)\right)`$, while the lower is constructed by assuming they were in the 3rd; $`\frac{1}{\pi }\left(p\left(0|\alpha \right)p\left(1|\alpha \right)p\left(n>1|\alpha \right)\right)`$. Of course neither one of these extremes can be the case as they are not physical. It is known that the Wigner distribution for any state must go to zero as $`x,p\mathrm{}`$. In addition, one can get a sense of the magnitude of the disagreement between $`W`$ and $`W^{}`$ by investigating the normalisation of the three curves; $`W^{}(x,p)𝑑x𝑑p=0.42`$, which is clearly not consistent with $`1`$. The same normalisation integral, however, gives values of $`10.32`$ and $`9.49`$ for the upper- and lower-bound curves, respectively. Clearly, the real Wigner function lies between $`W^{}`$ and the upper bound, but much closer to the former. Theoretically, one expects a significant additional contribution from the second excited state for $`|\alpha |1`$, but for larger values of $`|\alpha |`$ this should be cancelled out by higher populations in the third excited state. Near the origin, there was no population in states with $`n2`$, so there is no uncertainty.
Figure 4 shows a 3-D representation of the Wigner distribution measurement. Note the rotational symmetry, which indicates a lack of coherence between energy eigenstates. Usually processes such as the displacement operator create coherent superpositions of states such as in the case with our presented Husimi distribution. The lack of any coherence properties here is due to the combination of dephasing processes inherent in our optical latticemyrskog and the long delay between the population inversion creation and the time of measurement. Therefore the result of our state creation, as shown by figures 3 and 4 is an incoherent mixture of ground and 1st excited state atoms with most atoms being in the latter, as evidenced by the negative component of the Wigner distribution.
## IV Conclusion
Making use of a newly developed technique to measure state populations of the vibrational states of atoms in 1-$`\mu m`$ lattice wells, and our ability to perform arbitrary translations in phase space, we have reconstructed Husimi and Wigner phase space distributions for atoms in ground, coherent, and inverted states of oscillation in an optical lattice. This is the first complete quantum characterisation of the state of motion of atoms in such a system. The ground and coherent states are largely consistent with expectations, possessing essentially the same shape and width, although they indicate both some admixture of the excited state, which is understandable in light of our preparation procedure, and an apparent underestimate of the width of the ground-state wave function relative to the experimental measurement, which is probably due to inhomogeneities in the lattice beams. We observe a nonclassical signature in the Wigner function, reaching a maximum negative value of $`0.12`$ at the origin, in the case of the inverted distribution.
## V Acknowledgments
We would like to thank Matt Partlow for helpful discussions and assistance with the manuscript, and John Sipe for many informative conversations about Wigner distributions. We acknowledge financial support from NSERC, from the CIAR, and from the DARPA QuIST program (managed by the AFOSR under agreement no. F49620-01-1-0468).
Figure 1: Phase space representation of the Husimi distribution of a coherent state $`\left|\right|\beta |,\theta `$ for $`\left|\beta \right|=0.88`$, $`\theta =0.97`$. From the contour plot in the inset one can see the displacement of the Gaussian from the origin. Axis units are in $`x/2x_{rms}`$ where $`x_{rms}=96.3nm`$.
Figure 2: Cross sections of the Husimi distributions of the ground state (closed circles) and coherent state (open circles). As expected they have virtually the same width but the coherent state is displaced from the origin. In addition, the height and width, when compared to the true ground state of a harmonic oscillator of this frequency, show that there is some contamination of the state by first excited state atoms.
Figure 3: Cross section of the measured Wigner distribution. Negative values at the origin are indicative of a population inversion. Inset shows absolute upper (+) and lower ($`\times `$) bounds of the distribution as explained in the text. Axis units are in $`x/2x_{rms}`$ where $`x_{rms}=119.4nm`$
Figure 4: 3-D image of the measured Wigner distribution. Of note is the cylindrical symmetry showing a lack of coherence between the component energy eigenstates.
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# Deriving thermal lattice-Boltzmann models from the continuous Boltzmann equation: theoretical aspects
## Introduction
Following Lallemand & Luo, , the presently known lattice-Boltzmann equation (LBE) has not been able to handle realistic thermal (and fully compressible) fluids with satisfaction. Simulation of thermal lattice-Boltzmann equation is hampered by numerical instabilities when the local velocity increases. Readers are referred to this work for an excellent review of known works on thermal and compressible lattice-Boltzmann schemes.
Rigorously, fluid flow is never isothermal. Consider, for instance, a channel Poiseuille flow between two adiabatic solid surfaces. Due to the viscous conversion of mechanical in thermal energy, temperature will vary, attaining a minimum at the channel symmetry axis, where the local speed is a maximum. This temperature variation can be very small, but increases with the local average speed $`u`$ and with the fluid viscosity. If it is agreed that the temperature in a given site is related to the expected value E of the lattice-particles fluctuation kinetic energy, this temperature variation is drafted in LB athermal simulation, since E varies from site to site in accordance with the local macroscopic speed, $`u`$, attaining a minimum where $`u`$ is maximum. However, in athermal lattice-Boltzmann models, this thermal spatial non-equilibrium is not properly compensated by heat flow because athermal models were not conceived for correctly describing energy transfer. In this manner, since temperature gradients cannot be avoided in athermal LB simulation, in contrast with classical CFD isothermal simulation, they may become sources of numerical instability.
In conclusion, lattice-Boltzmann athermal equation, actually, deals with thermal problems and thermal and athermal lattice-Boltzmann models will be here considered using a single approach.
There are several features that let the lattice Boltzmann, regular-lattice based, framework far away from what it would desirable to be its starting point: the continuous Boltzmann equation. These features include the particles model, the collision model, the polynomial development used for the equilibrium distribution, the time discretization and the velocity discretization. Some of these main features are discussed in the following.
Collision model. More widely used lattice-Boltzmann collision models are based on a Bhatnagar, Gross and Krook (BGK) relaxation term, , firstly introduced in the lattice-Boltzmann framework by Qian et al., , and Chen et al., . Thermal lattice-Boltzmann schemes based on the BGK collision model use an increased number of discrete velocities and/or include higher order non-linear terms in the equilibrium distribution function (, , ), temperature dependent weights, , and temperature dependent velocities, . BGK single relaxation time collision term restricts the models to constant Prandtl number. The correct description of fluids and fluid flow requires multiple relaxation time models (MRT). A two-parameters model was introduced by He et al., , using two sets of distributions for the particles number density and for the thermodynamic internal energy, coupled through a viscous dissipation term. Full MRT models were firstly introduced in the LBE framework by d’Humières, , derived from the moments method, by making the moments and the velocity spaces isomorphic, . The main criticism to the moments method is that the highest order of the kinetic moments possible to be correctly described with the LBE equation is limited by the number of lattice velocities, , and high-order kinetic moments are not correctly described when all the b-moments in a b-discrete velocities set are considered. In currently produced works dealing with applications of the moments method, e.g. , the main worry is numerical stability and not the description of non-isothermal, multicomponent or immiscible fluids flows, which, effectively, require additional relaxation parameters with respect to BGK models. Dispersion equations are used as constraints for the adjustable parameters related to the short wave-length non-hydrodynamic moments and numerical stability is assured by buffering these higher frequency moments. This artificial shielding can be too dangerous for the complex flow structures that naturally appear when the Reynolds number increases. A correct description of the growing of these flow structures with the Reynolds number requires to increase the lattice dimensionality.
Time discretization. Most lattice-Boltzmann simulations are based on an explicit numerical scheme, with second order time-step, $`\delta `$, errors. Lattice BGK models, have been simulated with implicit numerical schemes, , , or LBE modified explicit numerical schemes, , with third order time step, O($`\delta ^3`$), errors. In spite of the fact that, in athermal models, this truncation error can be totally absorbed into the physical viscous term, in thermal models errors O($`\delta ^2`$) seriously affects the viscous heat dissipation term, .
Velocity discretization. Lattice-Boltzmann method is based on a finite set of discrete velocities $`\stackrel{}{c}_i`$ and weights $`\omega _i`$, judiciously chosen so as to ensure isotropy for the even-parity rank velocity tensors and, consequently, for the fluid transfer properties. He & Luo,, have directly derived some widely used lattices (D2Q9, D2Q6, D2Q7, D3Q27) from the continuous Boltzmann equation by discretization of the velocity space, using the Gauss-Hermite and Gauss -Radau quadrature of the Maxwellian $`\stackrel{}{u}`$-polynomial expansions. Unhappily, excluding the above mentioned lattices, the discrete velocity sets obtained by quadrature do not generate regular lattices. In this sense, Pavlo et al., , proposed a temperature dependent velocity model based on an octagonal lattice which is not space-filling but ensures the isotropy of 6<sup>th</sup>\- rank velocity tensors.
In this work, we present an attempt for deriving the lattice-Boltzmann equation, from the continuous Boltzmann equation trying to combine the following main features: multiple relaxation-times, O($`\delta ^3`$) time step errors and thermodynamic consistency in non-isothermal flow.
In contrast with the moments method, in present conception higher Reynolds flows (and non-isothermal flows) require to increase the number of the lattice discrete velocities and to increase the accuracy of the LBE equation with respect to its continuous counterpart.
The collision term $`\mathrm{\Omega }`$ in the linearized Boltzmann equation is modeled by expanding the distribution function $`f`$ in Hermite polynomial tensors $`\mathrm{\Psi }_\theta `$, which forms an orthogonal basis in the Hilbert space $``$ generated by $`h:𝒞^D`$, $`D`$ being the dimension of the velocity space. Considering that each term $``$ ($`\mathrm{\Psi }_\theta `$ ) is, itself, an element of $``$, this term is expanded as a linear combination of the same order-$`\theta `$ Hermite tensors through 2$`\theta `$-order relaxation tensors. Isotropy properties are used to reduce these tensors. The infinite series $`f^{eq}`$ ($`\varphi `$) is not truncated. Instead, after a chosen tensor order $`N`$, a Gross-Jackson procedure is used, , and the relaxation tensors are diagonalized.
It is shown that dH́umières moment equations are particular discrete forms of the derived model when the diagonalization constant is considered to be zero.
By performing a Chapman-Enskog analysis of the derived continuous model it is shown that the thermo-hydrodynamic macroscopic equations are correctly retrieved with a second-order model. Third-order models are only required for describing third-order coupling (Soret and Dufour effects) in multi-component systems ().
The derived kinetic model of the continuous Boltzmann equation is then discretized.
It is shown that an explicit numerical scheme with $`O`$($`\delta ^3`$) time step errors can be derived, using He at al. procedure, .
Velocity discretization is the most critical step in deriving lattice-Boltzmann equations.
For each $`N`$, the equilibrium distribution is taken as an $`n^{th}`$-degree Hermite development of the Maxwell-Boltzmann (MB) equilibrium distribution, in accordance with the constraints that are imposed by the physical problem.
Although, as it was shown in , velocity discretization of the most widely known lattices can be achieved by Gauss-Hermite and related quadratures, quadrature schemes did not succeed in generating multi-speed regular lattices, suitable for thermal problems, placing an, still, open question in the LBM framework. In present paper, it is shown that the integrating variable has an important role in defining the equilibrium distribution and the lattice-Boltzmann model: a) discretization based on the particles velocity $`\stackrel{}{c}`$, giving a set of discrete, constant, particle velocities $`\stackrel{}{c_i}`$, leads to temperature dependent weights $`\omega _i`$ (TDW), b) discretization based on the, temperature dependent, dimensionless velocity $`\stackrel{}{𝒞}`$ gives a set of temperature dependent particle velocities $`\stackrel{}{c_i}`$ (TDV).
In this context, it is shown that the thermal part $`g_i^{eq}`$ in He et al. two-distributions model, , can be formally retrieved from TDW models, as an $`\stackrel{}{u}`$-polynomial approximation with errors $`O\left(u\mathrm{\Theta }\right)`$ where $`\mathrm{\Theta }`$ is the temperature deviation.
Although a more complete theoretical analysis is still necessary, the consideration of TDV models appears to be suitable for thermal problems. However, in this case, the particles allocation rules, required for the local conservation of mass, momentum and energy, make the particles number-density, $`n`$, to be temperature dependent and the implicit temperature dependence of $`n`$ is difficult to manage, when performing a Chapman-Enskog analysis of the kinetic models. Finding the macroscopic behavior of these models is, still, in progress.
A simulation scheme suitable for non-isothermal problems is presented for the TDV model.
## 1 Boltzmann equation in the continuous velocity space
### 1.1 Development of the distribution function in Hermite polynomials
The Maxwell-Boltzmann equilibrium distribution, , can be written as an infinite series of Hermite polynomial tensors $`\mathrm{\Psi }_{\theta ,(r__\theta )}`$, ,
$$f^{eq}=\frac{e^{𝒞^2}}{\pi ^{D/2}}\left(\frac{m}{2kT}\right)^{D/2}\underset{\theta }{}a_{\theta ,(r__\theta )}^{eq}\mathrm{\Psi }_{\theta ,(r__\theta )},$$
(1)
where $`(r_\theta )`$ is a sequence of indexes $`r_1,r_2,\mathrm{}r_\theta `$ and repeated index means summation, $`\mathrm{\Psi }_0=1,\mathrm{\Psi }_{1,\alpha }=2𝒞_\alpha `$, $`\mathrm{\Psi }_{2,\alpha \beta }=2(𝒞_\alpha 𝒞_\beta \frac{1}{2}\delta _{\alpha \beta })`$, $`\mathrm{\Psi }_{3,\alpha \beta \gamma }=\frac{4}{3}(𝒞_\alpha 𝒞_\beta 𝒞_\gamma \frac{1}{2}\delta _{\alpha \beta }𝒞_\gamma \frac{1}{2}\delta _{\alpha \gamma }𝒞_\beta \frac{1}{2}\delta _{\beta \gamma }𝒞_\alpha )`$ and so on. The dimensionless particle velocity is $`\stackrel{}{𝒞}=\left(\frac{m}{2kT}\right)^{1/2}\stackrel{}{c}`$. These tensors are orthogonal in the Hilbert space $`,`$ satisfying
$$e^{𝒞^2}\mathrm{\Psi }_{\theta ,(r__\theta )}\mathrm{\Psi }_{\theta ,(s__\theta )}𝑑\stackrel{}{𝒞}=\lambda _\theta \mathrm{\Delta }_{(r__\theta )(s__\theta )},$$
(2)
where $`\mathrm{\Delta }_{(r__\theta )(s__\theta )}^{2\theta }`$ is a $`2\theta `$-order isotropic tensor, , and $`\lambda _\theta `$ is a constant. With $`\stackrel{}{𝒰}=\left(\frac{m}{2kT}\right)^{1/2}\stackrel{}{u}`$, the coefficients $`a_{\theta ,(r__\theta )}^{eq}`$ in Eq. (1) are the moments $`a__0^{eq}=n`$, $`a_{_{1,\alpha }}^{eq}=n𝒰_\alpha `$, $`a_{_{2,\alpha \beta }}^{eq}=n𝒰_\alpha 𝒰_\beta `$, $`a_{_{3,\alpha \beta \gamma }}^{eq}=n𝒰_\alpha 𝒰_\beta 𝒰_\gamma `$ and so on, which are dependents on the volumetric number of particles $`n`$, on the dimensionless macroscopic velocity $`\stackrel{}{𝒰}`$ and on the temperature $`T`$.
For each point $`\stackrel{}{x}`$ the distribution function $`\varphi `$ in the non-equilibrium part $`f^{neq}`$ $`=f^{eq}\varphi `$ can be developed in terms of the orthogonal basis $`\mathrm{\Psi }_{\theta ,(r__\theta )}`$, , , written in terms of the velocity fluctuation $`\stackrel{}{𝒞}_f=\frac{\stackrel{}{c}\stackrel{}{u}}{\left(\frac{2kT}{m}\right)^{1/2}}=\frac{\stackrel{}{C}}{\left(\frac{2kT}{m}\right)^{1/2}}=\stackrel{}{𝒞}\stackrel{}{𝒰}`$
$$\varphi =\underset{\theta }{}a_{\theta ,(r__\theta )}^\varphi (\stackrel{}{x},t)\mathrm{\Psi }_{\theta ,(r__\theta )}\left(\stackrel{}{𝒞}_f\right),$$
(3)
and coefficients $`a_\theta ^\varphi `$ can be related to the macroscopic moments of $`f`$. In this way, $`a_0^\varphi =0,`$ $`a_{1,\alpha }^\varphi =0`$. The coefficient $`a_{2,\alpha \beta }^\varphi `$ is related to the viscous stress tensor $`\tau _{\alpha \beta }`$ through
$$a_{2,\alpha \beta }^\varphi =\frac{\tau _{\alpha \beta }}{2P},$$
(4)
where $`P=nkT`$ is the thermodynamic pressure.
The fluctuation kinetic energy $`E(\stackrel{}{x},t)`$ is given by
$$E(\stackrel{}{x},t)=f\frac{1}{2}m\left(\stackrel{}{c}\stackrel{}{u}\right)^2𝑑\stackrel{}{c}=f^{eq}\frac{1}{2}m\left(\stackrel{}{c}\stackrel{}{u}\right)^2𝑑\stackrel{}{c.}$$
(5)
In this way
$$f^{neq}\frac{1}{2}m\left(C\right)^2𝑑\stackrel{}{C}=0,$$
(6)
or
$$f^{neq}\frac{1}{2}mC_\alpha C_\alpha 𝑑\stackrel{}{C}=\frac{1}{2}tr\left(\tau \right)=0.$$
(7)
In two-dimensions
$$\tau _{xx}+\tau _{yy}=0,$$
(8)
or
$$a_{2,xx}^\varphi +a_{2,yy}^\varphi =0.$$
(9)
For third-order moments
$`S_{\alpha \beta \gamma }`$ $`=`$ $`{\displaystyle fmc_\alpha c__\beta c_\gamma 𝑑\stackrel{}{c}}={\displaystyle f^{eq}mc_\alpha c__\beta c_\gamma 𝑑\stackrel{}{c}}+{\displaystyle f^{neq}mc_\alpha c__\beta c_\gamma 𝑑\stackrel{}{c}}`$ (10)
$`=`$ $`S_{\alpha \beta \gamma }^{eq}+S_{\alpha \beta \gamma }^{neq},`$
with
$$S_{\alpha \beta \gamma }^{eq}=\rho u_\alpha u_\beta u_\gamma +P\left(\delta _{\beta \gamma }u_\alpha +\delta _{\alpha \gamma }u_\beta +\delta _{\alpha \beta }u_\gamma \right).$$
(11)
For the non-equilibrium part,
$$S_{\alpha \beta \gamma }^{neq}=f^{neq}mc_{f\alpha }c_{f\beta }c_{f\gamma }𝑑\stackrel{}{c}+\left(\tau _{\beta \gamma }u_\alpha +\tau _{\alpha \gamma }u_\beta +\tau _{\alpha \beta }u_\gamma \right),$$
(12)
resulting, using $`a_{1,\alpha }^\varphi =0`$, the invariance property with respect to index permutation and Eq. (11):
$`P\left({\displaystyle \frac{2kT}{m}}\right)^{\frac{1}{2}}a_{3,\alpha \beta \gamma }^\varphi `$ $`=`$ $`{\displaystyle \frac{S_{\alpha \beta \gamma }}{2}}\left[\begin{array}{c}\frac{1}{2}\rho u_\alpha u_\beta u_\gamma +\frac{1}{2}P\left(\delta _{\beta \gamma }u_\alpha +\delta _{\alpha \gamma }u_\beta +\delta _{\alpha \beta }u_\gamma \right)\\ +\frac{1}{2}\left(\tau _{\beta \gamma }u_\alpha +\tau _{\alpha \gamma }u_\beta +\tau _{\alpha \beta }u_\gamma \right)\end{array}\right]`$ (15)
$``$ $`q_{\alpha \beta \gamma }.`$ (16)
When $`\beta `$ and $`\gamma `$ are contracted, defining $`ϵ_\alpha `$ to be the total energy flux along the direction $`\alpha `$,
$$P\left(\frac{2kT}{m}\right)^{\frac{1}{2}}a_{3,\alpha \beta \beta }^\varphi =ϵ_\alpha \left[\frac{1}{2}\rho u^2u_\alpha +P\left(\frac{D}{2}+1\right)u_\alpha +\tau _{\alpha \beta }u_\beta \right]=q_\alpha ,$$
(17)
where $`q_\alpha `$ is the net heat flux along the direction $`\alpha `$, i.e., the total energy flux $`ϵ_\alpha `$, subtracting from it, the flow of macroscopic kinetic energy $`\frac{1}{2}\rho u^2u_\alpha `$, the compression work $`P\left(\frac{D}{2}+1\right)u_\alpha `$ and the viscous work $`\tau _{\alpha \beta }u_\beta `$.
### 1.2 Collision term
Particles are supposed to be material points without volume and only able to exchange translational kinetic energy, but the collision term $`\mathrm{\Omega }`$ is, here, considered to take multiparticles collisions into account. Since, in this case, the collision term structure is not known, some assumptions are required. In this manner, near the equilibrium, $`\mathrm{\Omega }`$ is considered to be $`f^{eq}(\varphi )`$, the operator $``$ being a linear operator. This property was shown to be true for binary collisions, $`\left[14\right]`$ and is, here, extended for multiparticles collisions. When $`f`$ is near $`f^{eq}`$, Boltzmann equation reads
$$_tf+\stackrel{}{c}.f=\mathrm{\Omega }=f^{eq}(\varphi ).$$
(18)
Using the development, Eq. $`\left(\text{3}\right),`$
$$(\varphi )=\underset{\theta }{}a_{\theta ,(r_\theta )}^\varphi \left(\mathrm{\Psi }_{\theta ,(r_\theta )}\right).$$
(19)
The $`\theta `$-order tensor $`\left(\mathrm{\Psi }_{\theta ,(r_\theta )}\right)`$ is, itself, an element of the $`𝒞^D`$ space and can be developed in terms of the $`\theta `$-order Hermite tensors that belong to the orthogonal basis of this space,
$$\left(\mathrm{\Psi }_{_{\theta ,(r_\theta )}}\right)=\underset{(s_\theta )}{}\gamma _{_{(r_\theta ),(s_\theta )}}\mathrm{\Psi }_{_{\theta ,(s_\theta )}},$$
(20)
where $`\gamma _{(r_\theta ),(s_\theta )}`$ designate the $`(r_\theta ),(s_\theta )`$ components of $`2\theta `$-order relaxation tensors. Considering, as for binary collision, $``$ to be a self-adjoint operator, with non-positive eigenvalues,
$$\gamma _{_{(r_\theta ),(m_\theta )}}=\frac{e^{𝒞_f^2}\left(\mathrm{\Psi }_{_{\theta ,(r_\theta )}}\right)\mathrm{\Psi }_{_{\theta ,(m_\theta )}}𝑑\stackrel{}{𝒞}_f}{e^{𝒞_f^2}\left(\mathrm{\Psi }_{_{\theta ,(m_\theta )}}\right)^2𝑑\stackrel{}{𝒞}_f}0.$$
(21)
Using Einsteinś notation
$$(\varphi )=\underset{\theta }{}\gamma _{_{(r_\theta ),(s_\theta )}}a_{\theta ,(r_\theta )}^\varphi \mathrm{\Psi }_{_{\theta ,(s_\theta )}},$$
(22)
where repeated indexes mean summation.
Above equation is an infinite summation on $`\theta `$. When the terms above a chosen order N are diagonalised, following a Gross-Jackson procedure, ,
$$(\varphi )=\underset{\theta =0}{\overset{N}{}}\gamma _{_{(r_\theta ),(s_\theta )}}a_{\theta ,(r_\theta )}^\varphi \mathrm{\Psi }_{_{\theta ,(s_\theta )}}\gamma _{_{N+1}}\underset{\theta =N+1}{\overset{\mathrm{}}{}}\delta _{_{(r_\theta ),(s_\theta )}}a_{\theta ,(r_\theta )}^\varphi \mathrm{\Psi }_{_{\theta ,(s_\theta )}},$$
(23)
where
$$\delta _{_{(r_\theta ),(s_\theta )}}=\delta _{r__1s__1}\mathrm{}.\delta _{r__\theta s__\theta }.$$
(24)
In this way, using Eq. $`\left(\text{3}\right)`$
$$(\varphi )=\left[\underset{\theta =0}{\overset{N}{}}\lambda _{_{(r__\theta ),(s__\theta )}}a_{\theta ,(r__\theta )}^\varphi \mathrm{\Psi }_{_{\theta ,(s__\theta )}}\right]\gamma _{_{N+1}}\varphi ,$$
(25)
where $`\lambda _{_{(r__\theta ),(s__\theta )}}=\left(\gamma _{_{(r_\theta ),(s_\theta )}}+\gamma _{_{N+1}}\delta _{_{(r_\theta ),(s_\theta )}}\right)`$ is positive for all $`r_\theta ,s_\theta `$, since a) $`\lambda _{_{(r__\theta ),(s__\theta )}}=\gamma _{_{(r_\theta ),(s_\theta )}}`$ for all off-diagonal components and b) the diagonal components $`\gamma _{_{(r_\theta ),(r_\theta )}}`$ are negative with an absolute value that is greater than $`\gamma _{_{N+1}}`$ for all $`\theta `$ smaller or equal to $`N`$. Eq. (25) can be considered as an N<sup>th</sup>-order kinetic model to the collision term, with an absorption term $`\gamma __N\varphi `$ resulting from the diagonalization of the relaxation tensors after the given $`N`$. Therefore, all the moments of order higher than $`N`$ are collapsed into a single non-equilibrium term minimizing the truncation effects on the fine structure of the operator $``$ spectrum.
Although very little is known about the true collision term $`\mathrm{\Omega }`$ when multiple collisions are considered, Eq. (25) generates increasing accuracy models to $`\mathrm{\Omega }`$ when the distribution function $`f`$ is near the Maxwell-Boltzmann equilibrium distribution, $`f^{eq}`$. The only restrictions are: a) particles were considered as material points without volume and b) particles internal energy and long-range forces among the particles were not considered in the derivation.
When $`N=0`$ or $`N=1`$, Eq. (25) gives the well known BGK model, when all the collision operator spectra is replaced by a single relaxation term.
Each term in the sum, in Eq. (25), gives the relaxation to the equilibrium of second or higher order kinetic moments M<sub>θ</sub> that are not preserved in collisions, modulated by a $`\lambda _\theta `$ relaxation tensor. In Section 1.3, explicit expressions are given for the collision models. When the diagonalization constant is considered to be zero, i.e., when the series, Eq. (22), is truncated above N, replacing $`\mathrm{\Omega }=f^{eq}(\varphi )`$ in the Boltzmann equation, the inner products of the resulting equation by $`\mathrm{\Psi }_{_{\chi ,(s_\chi )}}`$ give
$$_ta_{\chi ,(r_\chi )}^f+\stackrel{}{c}.a_{\chi ,(r_\chi )}^f=\lambda _{(r_\chi )(s_\chi )}a_{\chi ,(s_\chi )}^{neq},$$
(26)
where the distribution function $`f`$ was developed following
$$f=\underset{\theta }{}e^{𝒞_f^2}a_{\theta ,(r_\theta )}^f\mathrm{\Psi }_{\theta ,(r_\theta )},$$
(27)
and
$$a_{\theta ,(r_\theta )}^{neq}=n\left(\frac{m}{2\pi kT}\right)^{D/2}a_{\theta ,(r_\theta )}^\varphi .$$
(28)
It can be easily seen that DH́umières moment equations (, ) are particular discrete forms of Eq. (26). Nevertheless, in dH́umières moments method, all the b-moments in a b-discrete velocities set are considered. It was shown in , that the number of degrees of freedom of a given lattice restricts the order $`n`$ of the kinetic moments with exact quadrature. This means that all the moments which order are greater than $`n`$ cannot be correctly described in this given lattice. As it was mentioned in the Introduction, in the moments method these high-frequency moments are forced to give consistent and numerical stable low-frequency macroscopic equations by using dispersion relations, decreasing the effect of numerical instability sources, but buffering the appearance of complex flow structures, when the Reynolds number increases.
### 1.3 Collision models for the continuous Boltzmann equation
In present section, the isotropy of 4<sup>th</sup> and 6<sup>th</sup> rank tensors will be used to give explicit forms for the second and third-order collision models, Eq. (25). Without any loss in the generality, we restrict ourselves to two-dimensional spaces.
#### 1.3.1 Second order model in the two-dimensional space
From Eq. (25)
$$\lambda _{_{(r__2),(s__2)}}a_{2,(r_2)}^\varphi \mathrm{\Psi }_{_{2,(s_2)}}=\lambda _{_{\alpha \beta \gamma \delta }}a_{2,\alpha \beta }^\varphi \mathrm{\Psi }_{_{2,\gamma \delta }}.$$
(29)
Requiring isotropy of 4<sup>th</sup> rank tensors and considering the symmetry with respect to index permutation,
$$\lambda _{_{\alpha \beta \gamma \delta }}=\lambda _\mu \left(\delta _{\alpha \beta }\delta _{\gamma \delta }+\delta _{\alpha \gamma }\delta _{\beta \delta }+\delta _{\alpha \delta }\delta _{\beta \gamma }\right).$$
(30)
In this way,
$`\lambda _{_{\alpha \beta \gamma \delta }}a_{2,\alpha \beta }^\varphi \mathrm{\Psi }_{_{2,\gamma \delta }}`$ $`=`$ $`\lambda _\mu \left[a_{2,\alpha \alpha }^\varphi \mathrm{\Psi }_{_{2,\gamma \gamma }}+a_{2,\alpha \beta }^\varphi \mathrm{\Psi }_{_{2,\alpha \beta }}+a_{2,\alpha \beta }^\varphi \mathrm{\Psi }_{_{2,\beta \alpha }}\right]`$ (33)
$`=`$ $`\lambda _\mu \left[\begin{array}{c}a_{2,xx}^\varphi \left(𝒞_{fx}^2\frac{1}{2}\right)+a_{2,yy}^\varphi \left(𝒞_{fy}^2\frac{1}{2}\right)+\\ 2a_{2,xy}^\varphi 𝒞_{fx}𝒞_{fy}\end{array}\right],`$
since $`a_{2,\alpha \alpha }^\varphi =0`$. Using Eq.(4)
$$\lambda _{_{\alpha \beta \gamma \delta }}a_{2,\alpha \beta }^\varphi \mathrm{\Psi }_{_{2,\gamma \delta }}=\frac{\lambda _\mu }{P}\left[\tau _{xx}\left(𝒞_{fx}^2\frac{1}{2}\right)+\tau _{yy}\left(𝒞_{fy}^2\frac{1}{2}\right)+2\tau _{xy}𝒞_{fx}𝒞_{fy}\right],$$
(34)
or, from Eq. (8), $`\tau _{xx}=\tau _{yy}`$
$$\lambda _{_{\alpha \beta \gamma \delta }}a_{2,\alpha \beta }^\varphi \mathrm{\Psi }_{_{2,\gamma \delta }}=\frac{\lambda _\mu }{P}\left[\tau _{xx}\left(𝒞_{fx}^2𝒞_{fy}^2\right)+2\tau _{xy}𝒞_{fx}𝒞_{fy}\right].$$
(35)
Second order model in two dimensions will be written as
$`^{(2)}(\varphi )`$ $`=`$ $`{\displaystyle \frac{\lambda _\mu }{P}}\left[\tau _{xx}\left(𝒞_{fx}^2{\displaystyle \frac{1}{2}}\right)+\tau _{yy}\left(𝒞_{fy}^2{\displaystyle \frac{1}{2}}\right)+2\tau _{xy}𝒞_{fx}𝒞_{fy}\right]`$ (36)
$`\gamma __3\varphi .`$
#### 1.3.2 <br>Third-order model
From Eq. (25)
$$\lambda _{_{(r__3),(s__3)}}a_{3,(r_3)}^\varphi \mathrm{\Psi }_{_{3,(s_3)}}=\lambda _{_{\alpha \beta \gamma \delta \zeta \eta }}a_{3,\alpha \beta \gamma }^\varphi \mathrm{\Psi }_{_{3,\delta \zeta \eta }}.$$
(37)
For isotropic fluids, tensor $`\lambda _{_{\alpha \beta \gamma \delta \zeta \eta }}`$ is a linear combination of five 6<sup>th</sup> order tensors given by the recurrence relation, ,
$$\mathrm{\Delta }_{r_1\mathrm{}.r_6}^{(6)}=\delta _{r_1r_j}\mathrm{\Delta }_{r_1\mathrm{}r_{j1},r_{j+1}\mathrm{}.r_4}^{(4)},$$
(38)
resulting
$$\lambda _{_{\alpha \beta \gamma \delta \zeta \eta }}a_{3,\alpha \beta \gamma }^\varphi \mathrm{\Psi }_{_{3,\delta \zeta \eta }}=\lambda __1a_{3,\alpha \beta \gamma }^\varphi \mathrm{\Psi }_{_{3,\alpha \beta \gamma }}+\lambda __2a_{3,\alpha \beta \beta }^\varphi \mathrm{\Psi }_{_{3,\alpha \beta \beta }}$$
(39)
### 1.4 Macroscopic thermohydrodynamic equations
Macroscopic thermohydrodynamic equations may be obtained from the Boltzmann equation by multiplying this equation by the mass m, the momentum m$`\stackrel{}{c}`$ and the kinetic energy$`\frac{1}{2}mc^2`$ of the particles and integrating the resulting equations in the $`\stackrel{}{c}`$ velocity space.
The mass conservation reads, as usually,
$$_t\rho +_\alpha \left(\rho u_\alpha \right)=0.$$
(40)
¿From the momentum preservation in collisions
$$_t\left(\rho u_\alpha \right)+_\alpha \left(\rho u_\alpha u_\beta +P\delta _{\alpha \beta }+\tau _{\alpha \beta }\right)=0,$$
(41)
where $`P`$ is the thermodynamics pressure, $`P=nkT`$ and $`\tau _{\alpha \beta }`$ is the viscous stress tensor.
When Eq.$`\left(\text{41}\right)`$ is multiplied by $`u_\alpha `$, the macroscopic kinetic energy, $`\frac{1}{2}\rho u^2`$, balance equation is obtained
$$_t\left(\frac{1}{2}\rho u_\alpha ^2\right)=P\mathrm{}.u+\tau _{\alpha \beta }_\beta u_\alpha _\beta \left(\frac{1}{2}\rho u_\alpha ^2u_\beta +P\delta _{\alpha \beta }u_\alpha +\tau _{\alpha \beta }u_\alpha \right).$$
(42)
The total energy conservation equation reads
$$_t\left(E+\frac{1}{2}\rho u_\alpha ^2\right)=_\beta \left[\left(\frac{1}{2}\rho u_\alpha ^2+E\right)u_\beta +\left(P\delta _{\alpha \beta }+\tau _{\alpha \beta }\right)u_\alpha +q_\beta \right],$$
(43)
where E is the thermodynamics internal energy
$$E=\frac{1}{2}m\left(\stackrel{}{c}\stackrel{}{u}\right)^2f𝑑\stackrel{}{c}=\frac{1}{2}m\left(\stackrel{}{c}\stackrel{}{u}\right)^2f^{eq}𝑑\stackrel{}{c}=\frac{D}{2}nkT.$$
(44)
The internal energy balance equation is obtained by subtracting Eq. (42) from Eq. (43),
$$_t\left(E\right)=(P\mathrm{}.u+\tau _{\alpha \beta }_\beta u_\alpha )_\beta [Eu_\beta +q_\beta ],$$
(45)
where $`(P\mathrm{}.u+\tau _{\alpha \beta }_\beta u_\alpha )`$ is the source term of internal energy.
Equations (40, 41 and 45) form a closed set of equations when the viscous stress tensor $`\tau _{\alpha \beta }`$ and the heat flux vector $`q_\beta `$ are known in terms of the spatial gradients of the first macroscopic moments, $`\rho `$, $`\stackrel{}{u}`$ and $`T`$ of the distribution function. This is accomplished when the Knudsen number, $`Kn0`$, by performing a Chapman-Enskog asymptotic analysis of the modelled Boltzmann equation.
### 1.5 Chapman Enskog analysis for the continuous model
Considering $`f^0`$ in the asymptotic expansion
$$f=f^0+Knf^1+\mathrm{},$$
(46)
to be the Maxwell-Boltzmann equilibrium distribution $`f^{eq}(n,\stackrel{}{u},T)`$, the zeroth order time derivative resulting from the Chapman-Enskog induced decomposition of the time derivative reads,
$`{\displaystyle \frac{1}{f^0}}{\displaystyle \frac{d_0f^0}{dt}}`$ $`=`$ $`2\left(𝒞_{f\alpha }𝒞_{f\beta }{\displaystyle \frac{1}{2}}\delta _{\alpha \beta }\right)_\beta u_\alpha {\displaystyle \frac{2}{D}}\left(𝒞_{f\alpha }^2{\displaystyle \frac{D}{2}}\right)\mathrm{}.\stackrel{}{u}`$ (47)
$`+\left({\displaystyle \frac{2kT}{m}}\right)^{1/2}\left(𝒞_f^2{\displaystyle \frac{D+2}{2}}\right)\stackrel{}{𝒞}_f.\mathrm{}\mathrm{ln}T.`$
#### 1.5.1 Second order model in two dimensions
Using Eqs. (47 and 36)
$`2\left(𝒞_{fx}^2{\displaystyle \frac{1}{2}}\right)_xu_x+2\left(𝒞_{fy}^2{\displaystyle \frac{1}{2}}\right)_yu_y+2𝒞_{fx}𝒞_{fy}\left(_xu_y+_yu_x\right)`$ (48)
$`\left[\left(𝒞_{fx}^2{\displaystyle \frac{1}{2}}\right)+\left(𝒞_{fy}^2{\displaystyle \frac{1}{2}}\right)\right]\mathrm{}.\stackrel{}{u}+\left({\displaystyle \frac{2kT}{m}}\right)^{1/2}\left(𝒞_f^22\right)\stackrel{}{𝒞}_f.\mathrm{}\mathrm{ln}T`$
$`=`$ $`{\displaystyle \frac{2\lambda _\mu }{P}}\left[\tau _{xx}\left(𝒞_{fx}^2{\displaystyle \frac{1}{2}}\right)+\tau _{yy}\left(𝒞_{fy}^2{\displaystyle \frac{1}{2}}\right)+2\tau _{xy}𝒞_{fx}𝒞_{fy}\right]\gamma __3\varphi .`$
For finding the correct expression of $`\tau _{\alpha \beta }`$, in terms of the spatial derivatives of the macroscopic variables, the inner product of the above equation by $`\left(𝒞_{f\alpha }𝒞_{f\beta }\frac{1}{2}\delta _{\alpha \beta }\right)`$ in the $`𝒞_{f\text{ }}`$ velocity space is performed. By multiplying the above equation by $`\left(𝒞_{fx}^2\frac{1}{2}\right)`$
$$\left(\lambda _\mu +\gamma __3\right)\frac{\tau _{xx}}{P}=_xu_x+\frac{1}{2}\mathrm{}.\stackrel{}{u}=\frac{1}{2}\left(_yu_y_xu_x\right).$$
(49)
Similarly
$$\left(\lambda _\mu +\gamma __3\right)\frac{\tau _{yy}}{P}=\frac{1}{2}\left(_xu_x_yu_y\right),$$
(50)
and
$$\left(\lambda _\mu +\gamma __3\right)\frac{\tau _{xy}}{P}=\frac{1}{2}\left(_xu_y+_yu_x\right).$$
(51)
These results give for the first and second viscosity coefficients,
$$\mu =\eta =\frac{nkT}{2\lambda _\mu +\gamma __3},$$
(52)
in the relation
$$\tau _{\alpha \beta }=\mu \left(_\alpha u_\beta +_\beta u_\alpha \right)+\eta \mathrm{}.\stackrel{}{u}.$$
(53)
Eq. (52) means that the first and the second viscosity coefficients are not independent quantities and this result must be considered as a limitation resulting from the continuous collision term itself, where the particles were considered as material points with translational degrees of freedom. In fact, this same result will be retrieved when using the third or higher order model for the collision term. The consideration of internal energy modes would be necessary for an up-grade of Eq. (52), .
The third-order moment $`a_{3,xxx}^\varphi `$ may be obtained by multiplying Eq. (48) by $`\left(𝒞_{fx}^2\frac{3}{2}\right)𝒞_{fx}`$. In this way,
$$a_{3,xxx}^\varphi =\frac{3}{4}\frac{\left(\frac{2kT}{m}\right)^{1/2}}{\gamma __3}_x\mathrm{ln}T.$$
(54)
Multiplying Eq. (48) by $`\left(𝒞_{fy}^2\frac{1}{2}\right)𝒞_{fx}`$,
$$a_{3,yyx}^\varphi =\frac{1}{4}\frac{\left(\frac{2kT}{m}\right)^{1/2}}{\gamma __3}_x\mathrm{ln}T.$$
(55)
But
$$q_x=P\left(\frac{2kT}{m}\right)^{1/2}\left(a_{3,xxx}^\varphi +a_{3,yyx}^\varphi \right)=\frac{P\left(\frac{2kT}{m}\right)}{\gamma __3T}_xT,$$
(56)
giving for the thermal conductivity
$$K=\frac{\left(D+2\right)nk^2T}{2m}\frac{1}{\gamma __3}.$$
(57)
In this manner, present second-order continuous kinetic model is thermodynamic consistent and able for analyzing non-isothermal and fully compressible flows. The thermal conductivity is related to $`\gamma __3`$ diagonalization constant. Consideration of third-order models will be, only, necessary in multi-component systems, for correctly describing third-order coupling: the Soret and Dufour effects, .
## 2 Discretization
In present sections an analysis is performed, trying to emphasize the theoretically identifiable effects of time and velocity discretization on the ability of the derived discrete models in retrieving the correct thermohydrodynamic equations, i.e., the full compressible Navier-Stokes equations and the thermodynamic internal energy balance equation, with the Fourier heat flux term.
### 2.1 Time discretization
Boltzmann equation with the kinetic model Eq. (25)becomes
$$\frac{d}{dt}f+\gamma __Nf=f^{eq}\underset{\theta =0}{\overset{N}{}}\lambda _{_{(r__\theta ),(s__\theta )}}a_{\theta ,(r__\theta )}^\varphi \mathrm{\Psi }_{_{\theta ,(s__\theta )}}+\gamma __Nf^{eq}.$$
(58)
For avoiding time-step errors $`𝒪\left(\delta ^2\right)`$, Boltzmann equation is integrated between $`t`$ and $`t+\delta `$, considering linear approximations for $`f^{eq}(\stackrel{}{x}+\stackrel{}{c}t^{},\stackrel{}{c},t+t^{})`$ and, also, for $`a_{\theta ,(r__\theta )}^{neq}`$, since $`a_{\theta ,(r__\theta )}^{neq}(\stackrel{}{x}+\stackrel{}{c}t^{},t+t^{})`$, when $`0t^{}\delta `$, . The result is
$`f(\stackrel{}{x}+\stackrel{}{c}\delta ,\stackrel{}{c},t+\delta )f(\stackrel{}{x},\stackrel{}{c},t)`$ (59)
$`=`$ $`\left(\gamma _N\delta \right){\displaystyle \frac{1}{2}}\left[f^{eq}(\stackrel{}{x}+\stackrel{}{c}\delta ,\stackrel{}{c},t+\delta )+f^{eq}(\stackrel{}{x},\stackrel{}{c},t)\right]`$
$`\left(\gamma _N\delta \right){\displaystyle \frac{1}{2}}\left[f(\stackrel{}{x}+\stackrel{}{c}\delta ,\stackrel{}{c},t+\delta )+f(\stackrel{}{x},\stackrel{}{c},t)\right]`$
$`+{\displaystyle \underset{\theta =0}{\overset{N}{}}}\left(\delta \lambda _{_{(r__\theta ),(s__\theta )}}\right)\mathrm{\Psi }_{_{\theta ,(s__\theta )}}f^{eq}(\stackrel{}{x},\stackrel{}{c},t)\text{ X}`$
$`\text{}{\displaystyle \frac{1}{2}}\left[a_{\theta ,(r__\theta )}^\varphi (\stackrel{}{x}+\stackrel{}{c}\delta ,t+\delta )+a_{\theta ,(r__\theta )}^\varphi (\stackrel{}{x},t)\right].`$
This corresponds to an implicit numerical scheme (in fact, a Crank-Nicholson scheme). Although implicit schemes are very easily manageable in the lattice-Boltzmann context, , , if one wants to avoid implicitness a new distribution can be defined as,
$$\stackrel{~}{f}=f+\gamma __N\frac{1}{2}\delta \left(ff^{eq}\right)\underset{\theta =0}{\overset{N}{}}\frac{\delta }{2}\lambda _{_{(r__\theta ),(s__\theta )}}a_{\theta ,(r__\theta )}^\varphi \mathrm{\Psi }_{_{\theta ,(s__\theta )}}f^{eq},$$
(60)
resulting
$`\stackrel{~}{f}(\stackrel{}{x}+\stackrel{}{c}\delta ,\stackrel{}{c},t+\delta )`$
$`=\stackrel{~}{f}(\stackrel{}{x},\stackrel{}{c},t)+`$ $`{\displaystyle \frac{\delta }{\tau __N+\frac{1}{2}\delta }}\left(f^{eq}\stackrel{~}{f}\right)+`$ (61)
$`{\displaystyle \frac{\delta }{\tau __N+\frac{1}{2}\delta }}{\displaystyle \underset{\theta =0}{\overset{N}{}}}\tau __N\lambda _{_{(r__\theta ),(s__\theta )}}a_{\theta ,(r__\theta )}^\varphi f^{eq}\mathrm{\Psi }_{_{\theta ,(s__\theta )}},`$
where $`\tau __N=\frac{1}{\gamma _N}`$. It must be observed that $`a_{\theta ,(r__\theta )}^\varphi `$ are the macroscopic moments of $`f`$ and not of $`\stackrel{~}{f}`$. Nevertheless, it can be shown from Eq. (60) that $`\stackrel{~}{a}_{\theta ,(s__\theta )}^\varphi `$ and $`a_{\theta ,(r__\theta )}^\varphi `$ are directly related by
$$\stackrel{~}{a}_{\theta ,(s__\theta )}^\varphi =\left(1+\frac{\delta }{2\tau __N}\right)a_{\theta ,(s__\theta )}^\varphi \frac{\delta }{2}\lambda _{_{(r__\theta ),(s__\theta )}}a_{\theta ,(r__\theta )}^\varphi ,$$
(62)
Consider, for instance, the second order model. In this case, it can be shown that
$`\lambda _{_{(r__\theta ),(s__\theta )}}a_{\theta ,(r__\theta )}^\varphi `$ $`=`$ $`\lambda _{_{\alpha \beta \gamma \delta }}a_{2,\alpha \beta }^\varphi =\lambda _\mu \left(\delta _{\alpha \beta }\delta _{\gamma \delta }+\delta _{\alpha \gamma }\delta _{\beta \delta }+\delta _{\alpha \delta }\delta _{\beta \gamma }\right)a_{2,\alpha \beta }^\varphi `$ (63)
$`=`$ $`2\lambda _\mu \left(a_{2,\gamma \delta }^{neq}\right),`$
since $`a_{2,\alpha \alpha }^{neq}=0`$ and $`\lambda _\mu `$ is required to be positive. Eq. (62) means
$$\stackrel{~}{a}_{2,\gamma \delta }^\varphi =\left(1+\frac{\delta }{2\tau __N}+\delta \lambda _\mu \right)a_{2,\gamma \delta }^\varphi .$$
(64)
### 2.2 Velocity discretization
Velocity discretization is the most critical step in presently proposed procedure. For athermal problems He and Luo, , have shown that some widely used sets of discrete velocities $`\left\{\stackrel{}{c}_i,i=1,\mathrm{}b\right\}`$ may be derived from the continuous velocity space by Gauss-Hermite (D2Q9, D3Q27) and Gauss-Radau (D2Q7) quadrature. All these sets are space-filling, in the sense that for every lattice site $`\stackrel{}{x}`$, $`\stackrel{}{x}+\stackrel{}{c}_i`$ points to another site in the lattice.
If it is agreed that quadrature is the bridge connecting the continuous and the discrete velocity space, discretization means to replace the entire continuous velocity space $`c^D`$ by some discrete velocities $`\stackrel{}{c}_i`$ satisfying the quadrature for all the kinetic moments of interest, i.e., for all the kinetic moments that are to be correctly described in lattice-Boltzmann simulation. Although it is highly desirable set $`\stackrel{}{c}_i`$ to be space-filling, this condition is not essential for the discretization itself.
When performing the quadrature, an integration variable must be chosen. If the dimensionless fluctuation velocity $`\stackrel{}{𝒞}_f=\frac{\stackrel{}{c}\stackrel{}{u}}{\left(\frac{2kT}{m}\right)^{1/2}}`$ is chosen as the integrating variable, considering $`\chi `$ to be a polynomial of degree $`r`$ in the velocity,
$$<\chi >=f\chi 𝑑\stackrel{}{c}=n\frac{1}{\pi ^{D/2}}e^{𝒞_f^2}\chi ^{}\left(\stackrel{}{𝒞}_f\right)𝑑\stackrel{}{𝒞}_f=n\underset{i=1}{\overset{b}{}}\omega _i\chi ^{}\left(\stackrel{}{𝒞}_{fi}\right),$$
(65)
where $`\stackrel{}{𝒞}_{fi}`$ is a discrete velocity (a constant vector), dependent, basically, on $`b`$ and on the kind of quadrature operation it is being performed, $`\chi ^{}\left(\stackrel{}{𝒞}_f\right)`$ is a polynomial in $`\stackrel{}{𝒞}_f`$ of degree $`r_p=r`$ $`+s`$ when it is related to a preserved moment and $`r_p=r+s+1`$, otherwise, , $`s`$ being the degree of the polynomial approximation to $`f^{eq}`$ and $`\omega _i`$ are the constant weights to be attributed to each discrete velocity $`\stackrel{}{𝒞}_{fi}`$. Exact quadrature restricts the highest value of $`r_p`$ to $`r_{p_m}`$. For a given class of quadrature we can write $`r_{p_m}=r_{p_m}(b)`$, in the sense that increasing $`b`$ enables higher degree polynomials to have exact quadrature.
For the first kinetic moment, $`n`$,
$$n=<1>=n\frac{1}{\pi ^{D/2}}e^{𝒞_f^2}1𝑑\stackrel{}{𝒞}_f=n\underset{i=1}{\overset{b}{}}\omega _i1=n\underset{i=1}{\overset{b}{}}\omega _i,$$
(66)
resulting,
$$f_i^{eq}=\omega _in.$$
(67)
This means that the discrete equilibrium distribution does not depend, explicitly, on the macroscopic velocity $`\stackrel{}{u}`$ and on the temperature $`T`$. These dependences are included in the particle velocities through,
$$\stackrel{}{c}_i=\stackrel{}{u}+\left(\frac{2kT}{m}\right)^{1/2}\stackrel{}{𝒞}_{fi}=\stackrel{}{c}_i(T,\stackrel{}{u}).$$
(68)
When the $`\stackrel{}{𝒞}_{fi}`$ generate a regular lattice, this choice is possible in LB framework, since LB simulation may, at least in principle, be performed in the ($`\stackrel{}{x},\stackrel{}{𝒞}_{fi})`$ space. Nevertheless, the physical grid ($`\stackrel{}{x},\stackrel{}{c}_i)`$, i.e., the physical grid points where the particles will be located after each time step, will be time dependent, simulation tends to be very cumbersome and, at a first sight, boundary conditions will be difficult to satisfy.
Another choice is the dimensionless particle velocity $`\stackrel{}{𝒞}=\frac{\stackrel{}{c}}{\left(\frac{2kT}{m}\right)^{1/2}}`$. This requires to rewrite the equilibrium distribution as in Eq. (1), but, now, the series
$$\underset{\theta }{}a_{\theta ,(r__\theta )}^{eq}(n,\stackrel{}{𝒰},T)\mathrm{\Psi }_{\theta ,(r__\theta )},$$
(69)
must be truncated somewhere. This is an important distinguishing point of discrete models, since in each continuous model presented in Section 1.3, although the full collision term is replaced by its N<sup>th</sup>-order approximation, the equilibrium distribution is, always, the full MB distribution.
Second-order approximations are widely used in athermal simulation, but thermohydrodynamics require third (or higher) order approximations for the equilibrium distribution,
$$f^{eq}=n\frac{e^{𝒞^2}}{\pi ^{D/2}}\left(\frac{m}{2kT}\right)^{D/2}\left[\begin{array}{c}1+2𝒞_\alpha 𝒰_\alpha +2\left(𝒞_\alpha 𝒞_\beta \frac{1}{2}\delta _{\alpha \beta }\right)𝒰_\alpha 𝒰_\beta +\\ \frac{4}{3}\left(𝒞_\alpha 𝒞_\beta \frac{3}{2}\delta _{\alpha \beta }\right)𝒞_\gamma 𝒰_\alpha 𝒰_\beta 𝒰_\gamma \end{array}\right],$$
(70)
which can be viewed as a third-degree polynomial expansion of the $`f^{eq}`$ dependence on $`\stackrel{}{𝒰}`$, with errors $`O(𝒰^4)`$.
After quadrature, the equilibrium distribution becomes
$$f_i^{eq}=\omega _in\left(\begin{array}{c}1+2𝒞_{i\alpha }𝒰_\alpha +2(𝒞_{i\alpha }𝒞_{i\beta }\frac{1}{2}\delta _{\alpha \beta })𝒰_\alpha 𝒰_\beta \\ +\frac{4}{3}\left(𝒞_{i\alpha }𝒞_{i\beta }\frac{3}{2}\delta _{\alpha \beta }\right)𝒞_{i\gamma }𝒰_\alpha 𝒰_\beta 𝒰_\gamma \end{array}\right),$$
(71)
where, as above, the weights $`\omega _i`$ and the velocity vectors $`\stackrel{}{𝒞}_i`$ are dependent on $`b`$ and on the kind of quadrature that was performed.
When $`\stackrel{}{u}=0`$, the equilibrium distribution is only dependent on the temperature $`T`$ through the number density of particles, $`n`$. Nevertheless, the particle velocities are temperature dependent,
$$\stackrel{}{c}_i=\left(\frac{2kT}{m}\right)^{1/2}\stackrel{}{𝒞}_i=\stackrel{}{c}_i\left(T\right).$$
(72)
A simulation alternative is presented, in this case, by redistributing the particles among adjacent sites, in accordance with allocation rules, locally preserving the mass, momentum and kinetic energy of the original packet. This strategy will be discussed in Section 2.3, leading to the establishment of temperature dependent velocity models (TDV).
Pavlo et al., , developed a TDV model based on an octagonal lattice, which is not space-filling but assures the isotropy of 6<sup>th</sup> rank tensors. It can be shown that this octagonal discrete velocities set can be retrieved using a Gauss-Radau quadrature, with 8 angular directions, giving $`r_{p_m}=7`$, instead of 5 as in the D2Q7 model and assuring the exact quadrature of third-order moments. Additional considerations can be found in Section 2.4.
Avoiding the $`\stackrel{}{c}_i`$ temperature dependence requires to consider the particles velocity $`\stackrel{}{c}`$ as the integrating variable when performing the quadrature, i.e., to let $`c^2`$ free from $`T`$ in the exponential part $`e^{𝒞^2}`$of the equilibrium distribution. This can be accomplished by writing
$$e^{\frac{\left(cu\right)^2}{\frac{2kT}{m}}}=\left(e^{𝒞_{fo}^2}\right)^{\frac{T_0}{T}},$$
(73)
where $`T_0`$ is a reference (and constant) temperature and $`\stackrel{}{𝒞}_{fo}=\frac{\stackrel{}{c}\stackrel{}{u}}{\left(\frac{2kT_0}{m}\right)^{1/2}}`$ is a new dimensionless fluctuation velocity referred to the temperature $`T_o`$. When $`T`$ is near $`T_0`$, i.e., when the departures from thermal equilibrium are small, the above expression may be developed in a Taylor series around $`\frac{T}{T_o}=1`$. Considering $`\mathrm{\Theta }=\frac{T}{T_o}1`$, this development gives
$$\left(e^{𝒞_{fo}^2}\right)^{\frac{T_0}{T}}=e^{𝒞_{fo}^2}\left[1+𝒞_{fo}^2\mathrm{\Theta }+\frac{1}{2}𝒞_{fo}^2\left(𝒞_{fo}^22\right)\mathrm{\Theta }^2+\mathrm{}\right],$$
(74)
which terms are increasing powers of $`𝒞_{fo}^2`$.
In this way, retaining just the first power in $`\theta `$,
$`f^{eq}`$ $`=`$ $`n\left({\displaystyle \frac{T_0}{T}}\right)^{D/2}\left[1+𝒞_{fo}^2\mathrm{\Theta }\right]{\displaystyle \frac{1}{\pi ^{D/2}}}\left({\displaystyle \frac{m}{2kT_0}}\right)^{D/2}`$ (77)
$`\times \text{ }e^{𝒞_o^2}\left[\begin{array}{c}1+2𝒞_{0,\alpha }𝒰_{0,\alpha }+2\left(𝒞_{0,\alpha }𝒞_{0,\beta }\frac{1}{2}\delta _{\alpha \beta }\right)𝒰_{0,\alpha }𝒰_{0,\beta }\\ +\frac{4}{3}\left(𝒞_{0\alpha }𝒞_{0\beta }\frac{3}{2}\delta _{\alpha \beta }\right)𝒞_{0\gamma }𝒰_{0\alpha }𝒰_{0\beta }𝒰_{0\gamma }\end{array}\right],`$
where $`\stackrel{}{𝒰}_0=\frac{\stackrel{}{u}}{\left(\frac{2kT_0}{m}\right)^{1/2}}`$.
In this case, the quadrature will give for the discrete equilibrium distribution,
$$f_i^{eq}=g_i(T,𝒞_{f0,i}^2)\omega _in\left[\begin{array}{c}1+2𝒞_{0,i\alpha }𝒰_{0,\alpha }+2\left(𝒞_{0,i\alpha }𝒞_{0,i\beta }\frac{1}{2}\delta _{\alpha \beta }\right)𝒰_{0,\alpha }𝒰_{0,\beta }\\ +\frac{4}{3}\left(𝒞_{i\alpha }𝒞_{i\beta }\frac{3}{2}\delta _{\alpha \beta }\right)𝒞_{i\gamma }𝒰_{0\alpha }𝒰_{0\beta }𝒰_{0\gamma }\end{array}\right],$$
(78)
where
$$g_i(T,𝒞_{f0,i}^2)=\left(\frac{T_0}{T}\right)^{D/2}\left[1+𝒞_{fo,i}^2\mathrm{\Theta }\right],$$
(79)
is a temperature dependent weight. When $`T>T_0,`$ $`g_i`$ reduces the amount of particles with zero velocity, redistributing them to the kinetic modes $`i`$ in accordance with $`𝒞_{f0,i}^2`$ and inversely when $`T<T_0`$. It is well known by LB practitioners that this redistribution is highly desirable in LB simulation and redistribution rules were empirically found by some authors (e.g., ).
Considering a linear approximation to the temperature non-equilibrium,
$$\left(\frac{T_0}{T}\right)^{D/2}1\frac{D}{2}\mathrm{\Theta }.$$
(80)
Eq. (79) can also be written as
$$g_i(T,𝒞_{f0,i}^2)=1\frac{D}{2}\mathrm{\Theta }+𝒞_{f0,i}^2\mathrm{\Theta }+O\left(\mathrm{\Theta }^2\right),$$
(81)
and the equilibrium distribution may be written as a sum of two distributions
$$f_i^{eq}=f_{i,n}^{eq}+f_{i,T}^{eq},$$
(82)
where, dropping-out the third-order term,
$$f_{i,T}^{eq}=\mathrm{\Theta }𝒞_{f0,i}^2\omega _in\left[1+2𝒞_{0,i\alpha }𝒰_{0,\alpha }+2\left(𝒞_{0,i\alpha }𝒞_{0,i\beta }\frac{1}{2}\delta _{\alpha \beta }\right)𝒰_{0,\alpha }𝒰_{0,\beta }\right].$$
(83)
This equilibrium distribution is related to the thermal distribution function $`g`$ in He et al. two-distributions model, . To fit the model into a D2Q9 lattice, He et al. have, further, replaced $`𝒞_{f0,i}^2`$ by $`\left(\stackrel{}{𝒞}_{0,i}\stackrel{}{𝒰}_0\right)^2`$, truncating all the terms $`O\left(𝒞_{0,i}^3\right)`$ and higher, after the multiplication. It can be easily seen that the resulting expression has second-order errors $`O\left(\mathrm{\Theta }𝒰_0\right)`$ limiting He et al.’s model to low local speeds.
We have preferred a somewhat different decomposition in Eq. (73), working with the particles velocity $`\stackrel{}{c}`$ and not with the fluctuation velocity $`\left(\stackrel{}{c}\stackrel{}{u}\right)`$, making
$$e^{\frac{c^2}{\frac{2kT}{m}}}=\left(e^{𝒞_o^2}\right)^{\frac{T_0}{T}},$$
(84)
resulting in a temperature dependent weights model (TDW), which equilibrium distribution is given by
$$f_i^{eq}=g_i(T,𝒞_{0,i}^2)\omega _in\left[\begin{array}{c}1+2\frac{𝒞_{0,i\alpha }𝒰_{0,\alpha }}{\frac{T}{T_0}}+2\left(𝒞_{0,i\alpha }𝒞_{0,i\beta }\frac{1}{2}\frac{T}{T_0}\delta _{\alpha \beta }\right)\frac{𝒰_{0,\alpha }𝒰_{0,\beta }}{\left(\frac{T}{T_0}\right)^2}\\ +\frac{4}{3}\left(𝒞_{0i\alpha }𝒞_{0i\beta }\frac{3}{2}\frac{T}{T_0}\delta _{\alpha \beta }\right)\frac{𝒞_{0i\gamma }𝒰_{0\alpha }𝒰_{0\beta }𝒰_{0\gamma }}{\left(\frac{T}{T_0}\right)^3}\end{array}\right],$$
(85)
where
$$g_i(T,𝒞_{0,i}^2)=1+\left(𝒞_{0,i}^2\frac{D}{2}\right)\mathrm{\Theta }.$$
(86)
Since $`\frac{T}{T_0},\left(\frac{T}{T_0}\right)^2,\mathrm{}`$ in Eq. (85), appears inside the polynomial expansion these terms must be, also, developed in $`\theta `$ for preserving consistency in the order of approximation. After multiplication, an expression for the equilibrium distribution in 2D problems can be written as
$$f_i^{eq}=\omega _in\left\{\begin{array}{c}1+\left(𝒞_{0,i}^21\right)\mathrm{\Theta }+2\left[1+\left(𝒞_{0,i}^22\right)\mathrm{\Theta }\right]\stackrel{}{𝒞}_{0,i}.\stackrel{}{𝒰}_0\\ +2\left(\stackrel{}{𝒞}_{0,i}\stackrel{}{.𝒰}_0\right)^2𝒰_0^2+\frac{4}{3}\left(\stackrel{}{𝒞}_{0,i}\stackrel{}{.𝒰}_0\right)^32(\stackrel{}{𝒞}_{0,i}.\stackrel{}{𝒰}_0)𝒰_0^2\end{array}\right\},$$
(87)
with errors $`O\left(\mathrm{\Theta }^2\text{}𝒰_0^2\mathrm{\Theta }\right)`$. A redistribution expression similar to the above one was obtained by Shan and He, . Large temperature deviations require to consider additional powers in $`\mathrm{\Theta }`$, in the Taylor expansion, Eq. (86) and to increase the number of discrete velocities in the lattice. This will be discussed in Section 2.4.3.
### 2.3 TDV model: allocation rules and collision step
Considering the analysis shown in the last section, TDV models appear to be a promissing alternative for LB simulation of non-isothermal problems, since no theoretical limitations were identified, related to the thermal non-equilibrium deviation $`\mathrm{\Theta }`$, as in the TDW models, Eq. (87), where the Taylor series in $`\mathrm{\Theta }`$, Eq. (86) was truncated after the first-order term. In present section this model is discussed, considering that temperature-dependent velocities require to modify the propagation step in the LB simulation scheme.
After quadrature, the third-order approximation to the equilibrium distribution Eq. (71) can be written in a general form as
$$\frac{f_k^{eq}}{\omega _kn}=\begin{array}{c}1+2a^2\frac{c_{k\alpha }^{}u_\alpha ^{}}{\frac{T}{T_0}}+2a^4\left(c_{k\alpha }^{}c_{k\beta }^{}\frac{1}{2a^2}\frac{T}{T_0}\delta _{\alpha \beta }\right)\frac{u_\alpha ^{}u_\beta ^{}}{\left(\frac{T}{T_0}\right)^2}\\ +\frac{4}{3}a^6\left(c_{k\alpha }^{}c_{k\beta }^{}\frac{3}{4a^2}\frac{T}{T_0}\delta _{\alpha \beta }\right)\frac{c_{k\gamma }^{}u_\alpha ^{}u_\beta ^{}u_\gamma ^{}}{\left(\frac{T}{T_0}\right)^3}\end{array},$$
(88)
where $`a`$ is a lattice constant related to the lattice symmetry. The $`\omega _k`$ are the lattice-weights, defining the inner product
$$f.g=\underset{k}{}\omega _kf_kg_k.$$
(89)
In general, the discrete velocities set generated by quadrature will be not space-filling requiring simulation strategies similar to the one proposed by Pavlo et al., . In space-filling lattices such as, e.g., the D2Q13H, the equilibrium distribution will keep the form given by Eq. (88) but the weights $`\omega _k`$ must be, empirically chosen for ensuring isotropy of even-parity rank velocity tensors up to the 6<sup>th</sup>-rank.
Particle velocities are given by
$$\stackrel{}{c}_k=\frac{h}{\delta }\stackrel{}{c}_k^{}.$$
(90)
where the lattice dimension $`h`$ is given by
$$\frac{h}{\delta }=a\left(\frac{2kT_0}{m}\right)^{1/2},$$
(91)
and
$$\stackrel{}{c}_k^{}=\sqrt{\frac{T}{T_0}}\stackrel{}{e}_k,$$
(92)
where $`\left\{\stackrel{}{e}_k,k=0,\mathrm{},b\right\}`$ are the lattice unity-vectors.
In present scheme, macroscopic velocity $`\stackrel{}{u}^{}=`$ $`\stackrel{}{u}/\left(\frac{h}{\delta }\right)`$, is calculated as $`n\stackrel{}{u}^{}=_kf_k\stackrel{}{c}_k^{}`$ where $`n=_kf_k`$. Temperature is found to be such that
$$n\frac{T}{T_0}=a^2\underset{k}{}f_k\left(\stackrel{}{c}_k^{}\stackrel{}{u}^{}\right)^2.$$
(93)
After discretization, the second-order model reads, in two-dimensions reads
$`\stackrel{~}{f_k}(\stackrel{}{x}+\stackrel{}{c}_k^{}\delta ,t+\delta )`$
$`=\stackrel{~}{f_k}(\stackrel{}{x},t)`$ $`+{\displaystyle \frac{\delta }{\tau __3+\frac{1}{2}\delta }}\left(f_k^{eq}\stackrel{~}{f}_k\right){\displaystyle \frac{\delta }{\tau __3+\frac{1}{2}\delta }}{\displaystyle \frac{1}{\left(1+\frac{\delta }{2\tau __3}+\frac{\delta }{2\tau __\mu }\right)}}`$
$`\times 6a^2`$ $`{\displaystyle \frac{\tau __3}{\tau __\mu }}{\displaystyle \frac{1}{\left(\frac{T}{T_0}\right)^2}}f_k^{eq}\left[\begin{array}{c}\stackrel{~}{\tau }_{xx}^{}\left(c_{fk,x}^2\frac{1}{2a^2}\frac{T}{T_0}\right)+\\ \stackrel{~}{\tau }_{yy}^{}\left(c_{fk,y}^2\frac{1}{2a^2}\frac{T}{T_0}\right)+2\stackrel{~}{\tau }_{xy}^{}c_{fk,x}^{}c_{fk,y}^{}\end{array}\right],`$ (96)
where $`\tau __3=1/\gamma __3`$ and $`\tau __\mu =1/\lambda __\mu `$.
The dimensionless viscous stress tensor is calculated as,
$$\tau _{\alpha \beta }^{}=\frac{\tau _{\alpha \beta }}{2a^2nkT_0}=\frac{1}{n}\underset{k}{}f_k^{neq}c_{k\alpha }^{}c_{k\beta }^{},$$
(97)
and the heat flux, as
$$q_\gamma ^{}=\frac{q_\gamma }{a^2kT_0\frac{h}{\delta }}=\underset{k}{}f_k^{neq}c_{fk\gamma }^{}c_{fk}^2.$$
(98)
The modified fluxes $`\stackrel{~}{\tau }_{\alpha \beta }^{}`$ and $`\stackrel{~}{q}_\gamma ^{}`$ needed in the simulation have similar expressions in terms of $`\stackrel{~}{f}_k^{neq}`$.
The particles that are present at site $`\stackrel{}{x}`$ after the collision step have a velocity $`\stackrel{}{c}_i^{}=\sqrt{\frac{T}{T_0}}\stackrel{}{e}_i`$. When $`\sqrt{\frac{T}{T_0}}1`$, the kinetic energy these particles have is just enough to enable them to jump to some intermediate position between $`\stackrel{}{x}`$ and $`\stackrel{}{x}+\stackrel{}{e}_i`$. In present discrete model these particles are redistributed to the next contiguous sites $`\stackrel{}{x}`$ and $`\stackrel{}{x}+\stackrel{}{e}_i`$ preserving the mass, kinetic energy and the momentum of the original particles packet $`\stackrel{~}{f}_i(\stackrel{}{x},t)`$. This is performed by using a leverś rule, making $`\stackrel{~}{f}_{i,0}(\stackrel{}{x},t+\delta )=\stackrel{~}{f}_i(\stackrel{}{x},t)\left(1\frac{T}{T_0}\right)`$ and $`\stackrel{~}{f}_{i,1}(\stackrel{}{x}+\stackrel{}{e}_i,t+\delta )=\stackrel{~}{f}_i(\stackrel{}{x},t)\left(\frac{T}{T_0}\right)`$, and by imposing to both amounts $`\stackrel{~}{f}_{i,0}`$ and $`\stackrel{~}{f}_{i,1}`$ the same non-integer velocity $`\stackrel{}{c}_i^{}`$ related to the particles-packet before the propagation. In this way model particles are velocity-memory particles in the propagation step. In the notation $`\stackrel{~}{f}_{ij}`$, the index $`i`$ is related to the site direction and $`j`$ means from what site these particles were come: $`j=0`$ means that the particles were come from the same site, $`j=1`$, from the site $`\stackrel{}{x}\stackrel{}{e}_i`$, $`j=2`$, from the site $`\stackrel{}{x}2\stackrel{}{e}_i`$ and so on.
In this manner, the particles propagated along the direction $`i`$ are redistributed among the departure site and the next-one contiguous site, but have their velocity unaltered, after the propagation. These rules preserve, locally, mass, momentum and energy and non-physical convection is avoided.
When $`1`$ $`\sqrt{\frac{T}{T_0}}2`$ , particles will be redistributed following
$$\stackrel{~}{f}_{i,1}(\stackrel{}{x}+\stackrel{}{e}_i,t+\delta )=\stackrel{~}{f}_i(\stackrel{}{x},t)\frac{1}{3}\left(4\frac{T}{T_0}\right),$$
(99)
and
$$\stackrel{~}{f}_{i,2}(\stackrel{}{x}+2\stackrel{}{e}_i,t+\delta )=\stackrel{~}{f}_i(\stackrel{}{x},t)\frac{1}{3}\left(\frac{T}{T_0}1\right).$$
(100)
Temperature ratios $`\sqrt{\frac{T}{T_0}}`$ greater than 2 are not expected in present work, since they would represent very strong thermal non-equilibrium.
After propagation, at time step $`t`$, all the distributions $`\stackrel{~}{f}_{i,j}(\stackrel{}{x},t)`$ are known in each site $`\stackrel{}{x}`$. The amounts $`\stackrel{~}{f}_{i,j}(\stackrel{}{x},t)`$ enable to calculate the equilibrium moments $`n(\stackrel{}{x},t)`$, $`u^{}(\stackrel{}{x},t)`$ and $`T(\stackrel{}{x},t)`$ and the equilibrium distribution, Eq. (88), with $`k=(i,j)`$.
The non-equilibrium distributions are given by $`\stackrel{~}{f}_{i,j}^{neq}=\stackrel{~}{f}_{i,j}f_{ij}^{eq}`$ and the viscous stress tensor can be thus calculated. In the collision step, collision will give a single distribution $`\stackrel{~}{f_i}^{}(\stackrel{}{x},t)`$ from the several $`j`$ amounts $`\stackrel{~}{f}_{i,j}`$ in the direction $`i`$ of the site $`\stackrel{}{x}`$ at time $`t`$,
$$\stackrel{~}{f_i}^{}(\stackrel{}{x},t)=\underset{j}{}\left\{\begin{array}{c}\stackrel{~}{f}_{i,j}+\frac{\delta }{\tau __N+\frac{1}{2}\delta }(f_{ij}^{eq}\stackrel{~}{f}_{i,j})\frac{\delta }{\tau __N+\frac{1}{2}\delta }\frac{6a^2\frac{\tau _3}{\tau _\mu }}{\left(1+\frac{\delta }{2\tau __N}+\frac{\delta }{2\tau _\mu }\right)}\times \\ \text{ }\frac{1}{\left(\frac{T}{T_0}\right)^2}f_{ij}^{eq}\left[\begin{array}{c}\stackrel{~}{\tau }_{xx}^{}\left(c_{f,ij,x}^2\frac{1}{2a^2}\frac{T}{T_0}\right)+\stackrel{~}{\tau }_{yy}^{}\left(c_{f,ij,y}^2\frac{1}{2a^2}\frac{T}{T_0}\right)\\ +2\stackrel{~}{\tau }_{xy}^{}c_{f,ij,x}^{}c_{f,ij,y}^{}\end{array}\right]\end{array}\right\}.$$
(101)
These particles will move to the next sites with the site velocity $`\stackrel{}{c^{}}_i=\sqrt{\frac{T}{T_0}}\stackrel{}{e}_i`$, where $`T=T(\stackrel{}{x},t)`$ is the site temperature, calculated using Eq.( 93), just after the propagation step and with $`k`$ replaced by $`(i,j)`$.
In the propagation step, a fraction of the low-speed particle-packets, in sites where $`T`$ is small, will remain in its departure site. The overall effect is to increase the number-density of particles in the sites where the temperature is small and to decrease it in the sites where the temperature is high, emulating the temperature dependence of the MB distribution. In this way, the equilibrium distribution is temperature dependent even when the macroscopic velocity $`\stackrel{}{u}`$ is zero, i.e., when $`f_i^{eq}=\omega _in`$, because $`n`$ is temperature dependent. Also, the temperature derivative of the equilibrium distribution will vary inversely with the temperature, similarly to its continuous counterpart. This is an important condition for retrieving the temperature derivative term in the second member of Eq.(47).
Nevertheless, the temperature dependence of $`n`$ is related to the particles allocation rules used and further theoretical analysis of presently proposed procedure leading to a correct Chapman-Enskog analysis of TDV model is, currently, under investigation.
### 2.4 The nine bits lattice
In this section, the TDW and TDV models are explicitly written for the D2Q9 lattice and their feasibility for simulating isothermal and non-isothermal problems in this lattice is discussed.
#### 2.4.1 <br>TDW model
The D2Q9 lattice is not suitable for thermal problems since: i) 6<sup>th</sup> rank tensors are not isotropic, ii) energy transfer $`q_{\alpha \beta \gamma }`$ is not correctly described in this lattice and iii) the lattice-dimensionality is too small. However, the use of thermal models in the D2Q9 lattice can, possibly, be shown to be a promising alternative for simulating near-isothermal problems, when the main interest is to increase numerical stability, with respect to athermal simulation. In fact, in its kinetic theory concept, temperature is varying from site to site, with the local velocity $`\stackrel{}{u}`$. When this variation is not considered, i.e., when the growing of temperature gradients are not annihilated by heat flow, they can, possibly, become instability sources in the numerical scheme. Actually, in contrast with classical CFD simulation of isothermal Navier-Stokes equations, lattice-Boltzmann simulation is never isothermal and the consideration of heat flow should be helpful in the simulations.
For the D2Q9 lattice, $`r_{p_m}(b)=5=2X31`$. Gauss-Hermite quadrature gives,
$$\stackrel{}{𝒞}_{0,i}=\sqrt{\frac{3}{2}}\stackrel{}{e}_i$$
(102)
where $`\stackrel{}{e}_i`$ = $`(0,0),(1,0),(0,1),(1,0),(0,1),(1,1),(1,1),`$ $`(1,1),(1,1).`$ The weights $`\omega _i`$ are found to be $`\omega __0=4/9,\omega _i=1/9,i=1,\mathrm{},4`$ for the main axes and $`\omega _i=1/36,i=5,\mathrm{},8`$, for the diagonals.
Using a dimensionless macroscopic velocity $`\stackrel{}{u}^{}=\frac{\stackrel{}{u}}{\frac{h}{\delta }}`$ and dropping-out the third-order term, since, as commented above, the exact description of energy transfer is out of thought in this lattice, the equilibrium distribution in the TDW model will read
$$f_i^{eq}=\omega _in\left[\begin{array}{c}1+\frac{3}{2}\left(e_i^2\frac{2}{3}\right)\mathrm{\Theta }+3\left[1+\frac{3}{2}\left(e_i^2\frac{4}{3}\right)\mathrm{\Theta }\right]\stackrel{}{e}_i.\stackrel{}{u}^{}\\ +\frac{9}{2}(\stackrel{}{e}_i.\stackrel{}{u}^{})^2\frac{3}{2}\left(u^{}\right)^2\end{array}\right].$$
(103)
It can be easily seen that the equilibrium distribution given by Eq. (103) satisfies
$$n=\underset{i}{}f_i^{eq},$$
(104)
$$n\stackrel{}{u}=\underset{i}{}f_i^{eq}\stackrel{}{c}_i,$$
(105)
$$P=nkT.$$
(106)
The internal energy
$$E=\underset{i}{}f_i^{eq}\frac{1}{2}m(\stackrel{}{c}_i\stackrel{}{u})^2,$$
(107)
the momentum flux
$$\mathrm{\Pi }_{\alpha \beta }^{eq}=\underset{i}{}f_i^{eq}mc_{i\alpha }c_{i\beta ,}$$
(108)
and the heat flux
$$q_\alpha ^{eq}=\underset{i}{}f_i^{eq}\frac{1}{2}m(\stackrel{}{c}_i\stackrel{}{u})^2\left(c_{i\alpha }u_\alpha \right),$$
(109)
are also retrieved as, respectively, $`\frac{D}{2}nkT`$, $`P\delta _{\alpha \beta }+\rho u_\alpha u_\beta `$ and $`0`$, with errors $`O(𝒰_0^2\mathrm{\Theta },\mathrm{\Theta }^2)`$.
#### 2.4.2 TDV model
All the restrictions given by Eqs.(104-106 and 107-109) are satisfied when the equilibrium distribution is given by Eq. (88), dropping-out the third-order term, i.e.,
$$f_k^{eq}=\omega _kn\left[1+3\frac{c_{k\alpha }^{}u_\alpha ^{}}{\frac{T}{T_0}}+\frac{9}{2}\left(c_{k\alpha }^{}c_{k\beta }^{}\frac{1}{3}\frac{T}{T_0}\delta _{\alpha \beta }\right)\frac{u_\alpha ^{}u_\beta ^{}}{\left(\frac{T}{T_0}\right)^2}\right].$$
(110)
Particles redistribution in TDV model should increase numerical stability with respect to the athermal simulation method, since, in this last scheme, particle allocation is not temperature dependent and sites with higher $`T`$ will rest with an excess amount in the number of particles, contributing to the enhancement of non-physical source-terms.
In Pavlo et al. , , TDV model, the use of an octagonal lattice, instead of the D2Q9 lattice, was considered for ensuring isotropy for the 6<sup>th</sup>-rank tensors. A second-order interpolation scheme is used for allocating the particles among adjacent sites in the propagation step, following a procedure that is similar to the one exposed in Section 2.3, but the model particles are not velocity-memory particles and the D2Q9 $`\stackrel{}{e}_i`$ lattice-velocities are attributed to the redistributed packets. This is performed preserving mass, momentum and energy. After the propagation, the number density, the momentum and the kinetic energy are calculated in each site. Previously to the collision step, the distributions $`f_i`$ are renormalized to account for the octagonal lattice velocities, satisfying the new local moment constraints. Successive transitions between the D2Q9 and the octagonal lattice can affect the isotropy properties of the 6<sup>th</sup>-rank tensors and, in our opinion, this is the only question that remains to be answered, since numerical viscosity effects can be avoided by reducing the spatial scale. In this manner, the problem is, nearly, the same as above, because the correct description of the heat flow vector cannot be assured in this model.
#### 2.4.3 Increasing the number of discrete velocities
After the very promising results of He and Luo’s work, , when the D2Q9, D2Q7 and D3Q27 lattices were found by exact quadrature, it is, now, apparent that the integer multiplicative factors that are required by regular-lattice velocity components are a too strong restriction: for their exact quadrature, Maxwellian-shaped curves require the roots $`c_{i\alpha }`$ to be progressively close when their absolute values increase.
For some regular lattices, the weights $`\omega _i`$ can be found to give even-parity isotropic tensors. In this manner, isotropy up to the 6<sup>th</sup>-rank tensors is assured for the D2Q13H lattice by choosing $`\omega _0=132/300`$, for the zero velocity, $`\omega _1=27/300`$ for the first 6 velocity vectors, with $`\left|\stackrel{}{e}_i\right|=1`$ and $`\omega _2=1/300`$ for the next 6 velocity vectors, with $`\left|\stackrel{}{e}_i\right|=\sqrt{3}`$. In this lattice $`h/\delta =\sqrt{\frac{10kT_0}{3m}\text{.}}`$
Simulation of non-isothermal problems is possible in this lattice by using third-order approximations to the MB equilibrium distribution, derived, alternatively, from Eq. (87) in the TDW model or from Eq. (88) in the TDV model.
In TDW models, improving the approximation, Eq. (87), to the MB equilibrium distribution, considering $`\mathrm{\Theta }^2`$ powers in Eq. (74), also requires to take 4<sup>th</sup>-degree powers of $`𝒞_0`$ into account and, consequently, to consider 4<sup>th</sup>\- degree polynomial expansions in $`𝒰_0`$ for the equilibrium distribution. This has been performed for a D2Q21S lattice, choosing the weights $`\omega _i`$ for giving even-parity isotropic tensors up to the 6<sup>th</sup>-rank and for retrieving the equilibrium conditions, Eqs.(104-106 and 107-109), with errors $`O(𝒰_0^3\mathrm{\Theta },\mathrm{\Theta }^3,𝒰_0^5,𝒰_0\mathrm{\Theta }^2,\mathrm{})`$. The D2Q21S is composed by the zero velocity and by the superposition of five square lattices with a $`45^0`$ shift between each two sequential lattices, with speeds $`0`$, $`1`$, $`\sqrt{2}`$, $`2`$, $`2\sqrt{2}`$ and $`3`$. The equilibrium distribution was found as
$$f_i^{eq}=\omega _in\left[\begin{array}{c}1+\left(e_i^21\right)\mathrm{\Theta }+\left(12e_i^2+e_i^4/2\right)\mathrm{\Theta }^2+2\left[1+\left(e_i^22\right)\mathrm{\Theta }\right]\stackrel{}{e_i}.\stackrel{}{u}^{}\\ +2(\stackrel{}{e_i}.\stackrel{}{u}^{})^2(1+(e_i^23)\mathrm{\Theta })\left(u^{}\right)^2(1+(e_i^22)\mathrm{\Theta })\\ +\frac{4}{3}(\stackrel{}{e}_i.\stackrel{}{u}^{})^32(\stackrel{}{e}_i.\stackrel{}{u}^{})\left(u^{}\right)^2+\frac{2}{3}(\stackrel{}{e}_i.\stackrel{}{u}^{})^4\\ 2(\stackrel{}{e}_i.\stackrel{}{u}^{})^2\left(u^{}\right)^2+\frac{1}{2}\left(u^{}\right)^4\end{array}\right].$$
(111)
with weights $`\omega _i=\{\frac{379}{1152},\frac{41}{384},\frac{5}{96},\frac{31}{3840},\frac{1}{1536},\frac{1}{5760}\}`$.
## 3 Conclusion
As a kinetic method, LBM is based on a discrete and finite approximation to the continuous Boltzmann equation giving an approximated solution to the particles distribution function. In this manner, all the macroscopic moments of interest will be affected by the accuracy of the modelled LBE when compared to its full continuous counterpart.
MRT increasing accuracy models to the collision term $`\mathrm{\Omega }`$ were derived collapsing the higher-order terms, related to the relaxation of the higher-order moments, into a single non-equilibrium term, minimizing the truncation effects on the fine structure of the collision operator spectrum.
Thus, in contrast with the moments method, only the hydrodynamic moments that are required to be correctly described, are considered in the present models. This strategy avoids the use of dispersion relations for buffering undesirable effects of the high-frequency, non-hydrodynamic moments, but claims to increase the lattice dimensionality.
Time discretization requires implicit or modified explicit numerical schemes for avoiding $`𝒪(\delta ^2)`$ time step errors and non-physical source terms in the internal energy balance equation.
For the discretization of the velocity space, the requirements for the Boltzmann equation quadrature to be exact establishes the minimal discrete velocity set that is needed for the lattice, in accordance with the order of approximation for the kinetic model that is intended to be used. Nevertheless, these sets do not generate regular lattices when multi-speed models appropriated for thermal problems are considered.
When performing the quadrature, it was shown that the integrating variable has an important role in defining the equilibrium distribution and the lattice-Boltzmann model, leading, alternatively, to temperature dependent velocities (TDV) and to temperature dependent weights (TDW) lattice-Boltzmann models.
In TDV models no theoretical limitations were identified, related to the thermal non-equilibrium deviation $`\mathrm{\Theta }`$, as in the TDW models.
Finding the macroscopic behavior of TDV models through a rigorous Chapman-Enskog analysis is difficult due to the implicit number density of particles dependence on the temperature and is, still, in progress.
Acknowledgements
Authors are greatly indebted to ANP (Brazilian Petroleum Agengy), CNPq (Brazilian Research Council), Finep (Brazilian Agency for Research and Projects) and Petrobras (Brazilian Petroleum Company) for the financial support.
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# A Note on Perturbations During a Regular Bounce
## Abstract
We point out an inconsistency in a method used in the literature for studying adiabatic scalar perturbations in a regular bouncing universe (in four dimensions). The method under scrutiny consists of splitting the Bardeen potential into two pieces with independent evolutions, in order to avoid a singular behavior at the boundaries of the region where the null energy condition (NEC) is violated. However, we argue that this method violates energy-momentum conservation.
We then introduce a novel method which provides two independent solutions for the Bardeen potential around the boundaries, even in the case of adiabatic perturbations. The two solutions are well behaved and not divergent.
preprint: BROWN-HET-1449
In recent years, regular toy-models of a bouncing universe have received a lot of attention due to string inspired singular bouncing models such as the cyclic scenario. Of particular interest is the evolution of scalar perturbations, since their spectrum is directly observable in the anisotropies of the cosmic background radiation and the large scale structure of the universe. The majority of toy models are four dimensional, requiring at least two matter fields, of which at least one has to violate the null energy condition (NEC).
If one focuses on adiabatic perturbations, the evolution equation for the Bardeen potential becomes singular at the boundaries of the NEC violating region before and after the bounce. One way out of this dilemma was advocated in Peter:2002cn and subsequently used e.g. in Peter:2003 ; Finelli:2003mc ; Battefeld:2005cj : One splits the Bardeen potential in two components, each of which satisfies a regular second order differential equation.
In this note we first show that this method is inconsistent because the fluid conservation equations are violated. Therefore all models using this method have to be re-evaluated. We then introduce a novel method, providing two independent and well behaved solutions for the Bardeen potential around the boundaries of the region with NEC violation, even in the case of adiabatic perturbations.
For simplicity, we work with a two fluid model with $`T_\nu ^\mu =T_{(a)\nu }^\mu \pm T_{(b)\nu }^\mu `$ and equations of state $`p_{(l)}=w_l\rho _{(l)}`$ for $`l=a,b`$. Note that a bounce for the scale factor $`a`$ occurs only in case of a negative sign in front of $`T_{(b)\nu }^\mu `$ <sup>3</sup><sup>3</sup>3$`w_b=2w_a+1`$ corresponds to the bounce in Battefeld:2005cj and $`w_b=2w_a+1/3`$ corresponds to the one in Finelli:2003mc .. Perturbing the metric
$`ds^2`$ $`=`$ $`a^2\left[(1+2\mathrm{\Phi })d\eta ^2(12\mathrm{\Phi })\delta _{ij}dx^idx^j\right],`$ (1)
where $`\mathrm{\Phi }`$ is the Bardeen potential (longitudinal gauge, no anisotropic stress, see Mukhanov:1990me for details), and perturbing also the energy momentum tensors
$`(\delta T_{(l)\nu }^\mu )`$ $`=`$ $`\rho _{(l)}\left(\begin{array}{cc}\xi _{(l)}& (1+w_l)V_{(l),i}\\ (1+w_l)V_{(l),i}& (w_l\xi _{(l)})\delta _j^i\end{array}\right)`$ (4)
where $`\delta p_{(l)}=w_l\delta \rho _{(l)}`$ and $`\xi _{(l)}:=\delta \rho _{(l)}/\rho _{(l)}`$, the perturbed Einstein equations for $`𝒦=0`$ (a spatially flat universe) read
$`^2\mathrm{\Phi }3(\mathrm{\Phi }+\mathrm{\Phi }^{})={\displaystyle \frac{a^2}{2}}\kappa ^2\left(\rho _{(a)}\xi _{(a)}\pm \rho _{(b)}\xi _{(b)}\right),`$
(5)
$`\mathrm{\Phi }^{\prime \prime }+3\mathrm{\Phi }^{}+(2^{}+^2)\mathrm{\Phi }={\displaystyle \frac{a^2}{2}}\kappa ^2(w_a\xi _{(a)}\rho _{(a)}`$
$`\pm w_b\xi _{(b)}\rho _{(b)}),`$ (6)
$`[\mathrm{\Phi }+\mathrm{\Phi }^{}]_{,i}={\displaystyle \frac{a^2}{2}}\kappa ^2(\rho _{(a)}V_{(a),i}(1+w_a)`$
$`\pm \rho _{(b)}V_{(b),i}(1+w_b)),`$ (7)
with $`\kappa ^2=8\pi /M_p^2`$. The energy conservation equations for each fluid, assuming no non-gravitational interactions between the fluids, are
$`\rho _{(l)}\left((1+w_l)\left[^2V_{(l)}3\mathrm{\Phi }^{}\right]+\xi _{(l)}^{}\right)=0.`$ (8)
In the long wavelength limit the above relation yields
$`{\displaystyle \frac{\xi _{(l)}^{}}{1+w_l}}=3\mathrm{\Phi }^{},`$ (9)
which is another way of stating entropy conservation, that is
$`S^{}={\displaystyle \frac{\xi _{(a)}^{}}{1+w_a}}{\displaystyle \frac{\xi _{(b)}^{}}{1+w_b}}=0.`$ (10)
One can simply impose the adiabaticity condition by setting initially $`S(\eta _{in}):=0`$ and thus <sup>4</sup><sup>4</sup>4While one can not use the energy conservation equations for separate fluids in the case discussed in Battefeld:2005cj to derive this condition, one can still arrive at the same relation due to the specific form of the corrections to the energy momentum tensor in Battefeld:2005cj .
$`{\displaystyle \frac{\xi _{(a)}}{1+w_a}}={\displaystyle \frac{\xi _{(b)}}{1+w_b}}`$ (11)
has to hold.
Substituting this result back into (5) and (6) and then combining the two, one arrives at a second order equation for $`\mathrm{\Phi }`$:
$`0=\mathrm{\Phi }^{\prime \prime }\left[\rho _a(1+w_a)\pm \rho _b(1+w_b)\right]`$ (12)
$`+3\mathrm{\Phi }^{}\left[\rho _a(1+w_a)^2\pm \rho _b(1+w_b)^2\right]`$
$`+[(w_a(w_a+1)\rho _a\pm w_b(w_b+1)\rho _b)^2`$
$`+2^{}\left(\rho _a(1+w_a)\pm \rho _b(1+w_b)\right)`$
$`+^2(\rho _a(1+w_a)(1+3w_a)\pm \rho _b(1+w_b)(1+3w_b))]\mathrm{\Phi }.`$
In case of a minus sign, which is needed for a bounce to occur, or some negative $`w_i`$ this equation becomes singular at the boundaries of the NEC violating region (see also Finelli:2003mc ).
In the appendix of Peter:2002cn a split of $`\mathrm{\Phi }`$ into $`\mathrm{\Phi }_a+\mathrm{\Phi }_b`$ is suggested, such that each $`\mathrm{\Phi }_l`$ satisfies
$`^2\mathrm{\Phi }_l3(\mathrm{\Phi }_l+\mathrm{\Phi }_l^{})={\displaystyle \frac{a^2}{2}}\kappa ^2\rho _{(l)}\xi _{(l)},`$ (13)
$`\mathrm{\Phi }_l^{\prime \prime }+3\mathrm{\Phi }_l^{}+(2^{}+^2)\mathrm{\Phi }_l={\displaystyle \frac{a^2}{2}}\kappa ^2w_l\xi _{(l)}\rho _{(l)}.`$ (14)
By adding these equations one arrives at the original Einstein equations (5) and (6). Combining (13) and (14) one can derive for each $`\mathrm{\Phi }_l`$ the equation
$`0`$ $`=`$ $`\mathrm{\Phi }_l^{\prime \prime }+3(1+w_l)\mathrm{\Phi }_l^{}`$ (15)
$`+\left(w_l^2+2^{}+(1+3w_l)^2\right)\mathrm{\Phi }_l.`$
Note that these equations are regular for each $`\mathrm{\Phi }_l`$. This is already the first hint that the method is doubtful, since the singular behavior vanished miraculously. Another reason to doubt this method is the fact that by just using equations (5) and (6), with no extra constraint like conservation of energy for each fluid or adiabaticity, one seems to be able to calculate the evolution of each $`\mathrm{\Phi }_l`$ and subsequently $`\mathrm{\Phi }`$ itself. However, it is obvious that in the equations (5) and (6) three unknowns appear, so that two equations are insufficient to calculate their evolution. Therefore one can not expect the solutions of (13) and (14) to be consistent with the conservation of the energy-momentum tensor or other constraint equations.
One easy way to see this inconsistency is to look at the solutions in the long wavelength limit where equation (13) simplifies to
$`3(\mathrm{\Phi }_l+\mathrm{\Phi }_l^{})={\displaystyle \frac{a^2}{2}}\kappa ^2\rho _{(l)}\xi _{(l)}.`$ (16)
Adding the time derivative of the above equation to $`\times `$ itself and $`3\times `$ (14) yields
$`\rho _{(a)}\xi _{(a)}^{}`$ $`=`$ $`3\mathrm{\Phi }_a^{}\left(\rho _{(a)}(1+w_a)+\rho _{(b)}(1+w_b)\right),`$ (17)
$`\rho _{(b)}\xi _{(b)}^{}`$ $`=`$ $`3\mathrm{\Phi }_b^{}\left(\rho _{(b)}(1+w_b)+\rho _{(a)}(1+w_a)\right),`$ (18)
where we also used the background Einstein equations and $`\rho _i^{}=3(1+w_i)\rho _i`$. This result, together with the conservation equation (9), implies that the solutions of (15) have to satisfy
$`{\displaystyle \frac{\mathrm{\Phi }_a^{}}{\rho _{(a)}(1+w_a)}}={\displaystyle \frac{\mathrm{\Phi }_b^{}}{\rho _{(b)}(1+w_b)}}.`$ (19)
One might think that this constraint is satisfied in the case of adiabatic perturbations, but it is not as we shall see now. Using (11) in (16) yields
$`{\displaystyle \frac{\mathrm{\Phi }_a+\mathrm{\Phi }_a^{}}{\rho _{(a)}(1+w_a)}}={\displaystyle \frac{\mathrm{\Phi }_b+\mathrm{\Phi }_b^{}}{\rho _{(b)}(1+w_b)}},`$ (20)
which can further be simplified by using (19) to
$`{\displaystyle \frac{\mathrm{\Phi }_a}{\rho _{(a)}(1+w_a)}}={\displaystyle \frac{\mathrm{\Phi }_b}{\rho _{(b)}(1+w_b)}}.`$ (21)
This together with (19) implies
$`{\displaystyle \frac{\mathrm{\Phi }_a^{}}{\mathrm{\Phi }_a}}={\displaystyle \frac{\mathrm{\Phi }_b^{}}{\mathrm{\Phi }_b}},`$ (22)
yielding $`\mathrm{\Phi }_a\mathrm{\Phi }_b`$. However, this is clearly in contradiction to (21), since the densities have a different dependency on conformal time if $`w_aw_b`$. Therefore the splitting method itself is inconsistent and can not be trusted to regularize (12).
We propose another mathematical technique that can be used instead to approximate the solutions of (12) in the vicinity of $`\rho _{tot}+p_{tot}=0`$ at $`\eta _{nec}`$. Equation (12) is a second order differential equation that in fourier space (suppressing the subscript on $`\mathrm{\Phi }_k`$), has the following general form:
$`A(\eta )\mathrm{\Phi }^{\prime \prime }+B(\eta )\mathrm{\Phi }^{}+C(k,\eta )\mathrm{\Phi }=0,`$ (23)
where $`A`$ and $`B`$ are related via
$`B=A^{},`$ (24)
since energy conservation for each fluid requires <sup>5</sup><sup>5</sup>5Again, although the energy conservation argument does not hold for Battefeld:2005cj this relation still remains valid.
$`\rho _l^{}=3\rho __l(1+w_l).`$ (25)
Our first goal is to derive one of the solutions around $`\eta _{nec}`$ perturbatively. We can then obtain the second solution by means of the Wronskian method.
Since $`a(\eta _{nec})0`$ we can Taylor expand $`\rho _l`$ and consequently $`A,B`$ and $`C`$ around $`\eta _{nec}`$:
$`A(\delta )=A_1\delta +A_2\delta ^2+A_3\delta ^3+\mathrm{},`$ (26)
$`B(\delta )=A_12A_2\delta 3A_3\delta ^2+\mathrm{},`$ (27)
$`C(\delta )=C_0+C_1\delta +C_2\delta ^2+\mathrm{},`$ (28)
where we defined $`\delta =\eta \eta _{nec}`$ and used $`A(\eta _{nec})=0`$. We assume that $`A_1`$, corresponding to the linear term of $`A(\delta )`$, does not vanish. This is usually the case and plays a crucial role in our analysis.
Furthermore, we assume the existence of an analytic solution around $`\eta _{nec}`$ (we will see bellow that it is a justified assumption <sup>6</sup><sup>6</sup>6Arguments in favor of an analytic $`\mathrm{\Phi }`$ where given in e.g. Peter:2003 or Bozza:2005xs , without relying on the splitting method. ), so that it can be written as
$`\mathrm{\Phi }_1=\alpha _0+\alpha _1\delta +\alpha _2\delta ^2+\mathrm{}.`$ (29)
Substituting relations (26)-(29) into (23) we obtain the following relations from the zeroth and first order equations in $`\delta `$
$`𝒪(\delta ^0)`$ $`0=A_1\alpha _1+C_0\alpha _0,`$ (30)
$`𝒪(\delta ^1)`$ $`0=(2A_2+C_0)\alpha _1+C_1\alpha _0.`$ (31)
In general, this can only be satisfied if $`\alpha _0=\alpha _1=0`$. Fortunately, this does not imply $`\mathrm{\Phi }_10`$, because the equations of higher order in $`\delta `$ can all be satisfied recursively. In fact, we can compute a complete power series solution for $`\mathrm{\Phi }_1`$
$`𝒪(\delta ^2)`$ $``$ $`\alpha _3={\displaystyle \frac{2A_2C_0}{3A_1}}\alpha _2`$
$`\mathrm{}`$
$`𝒪(\delta ^n)`$ $``$ $`\alpha _{n+1}={\displaystyle \frac{\underset{i=2}{\overset{n}{}}[i(3+n2i)A_{n+2i}C_{ni}]\alpha _i}{(n1)(n+1)A_1}}`$
with the consequence
$`\mathrm{\Phi }_1=\delta ^2+{\displaystyle \frac{2A_2C_0}{3A_1}}\delta ^3+\mathrm{}.`$ (33)
Knowing one of the solutions of (23), $`\mathrm{\Phi }_1`$, we can easily obtain the other solution, $`\mathrm{\Phi }_2`$, by using the Wronskian technique. The Wronskian for a second order differential equations is defined as
$`W=\mathrm{\Phi }_2^{}\mathrm{\Phi }_1\mathrm{\Phi }_1^{}\mathrm{\Phi }_2,`$ (34)
where $`\mathrm{\Phi }_1`$ and $`\mathrm{\Phi }_2`$ are the independent solutions of (23). Henceforth, one can calculate $`\mathrm{\Phi }_2`$ in terms of $`W`$ and $`\mathrm{\Phi }_1`$:
$`\mathrm{\Phi }_2\mathrm{\Phi }_1{\displaystyle \frac{W}{\mathrm{\Phi }_1^2}𝑑\eta }.`$ (35)
$`W`$ itself satisfies the first order differential equation
$`A(\eta )W^{}+B(\eta )W=0.`$ (36)
By substituting $`B(\eta )`$ from (24) we can solve the above equation for $`W`$ to
$`W(\eta )=\beta A(\eta ),`$ (37)
where $`\beta `$ is just a constant. Combining this result for $`W(\eta )`$ with (35) and our solution for $`\mathrm{\Phi }_1`$ from (33), we end up with
$`\mathrm{\Phi }_2={\displaystyle \frac{A_1}{2}}{\displaystyle \frac{8A_2C_0}{6}}\delta +A_3\delta ^2\mathrm{ln}(|\delta |)+O(\delta ^2).`$ (38)
Note that although this solution is not analytic at $`\eta =\eta _{nec}`$ ($`\delta =0`$), it is well behaved in the sense that both, the solution and its first derivative, remain continuous and finite. Thus, we have constructed the approximate form of the two independent solutions of (23) or subsequently (12). These can be used to match the solution on different sides of $`\eta _{nec}`$.
To summarize, we have shown explicitly that the splitting method, first introduced in Peter:2002cn and used to show the regularity of the Bardeen potential, is intrinsically inconsistent. All models using this method have to be re-evaluated <sup>7</sup><sup>7</sup>7The model in Battefeld:2005cj V.2 by the authors of this critique will be corrected shortly in the upcoming revision V.3. The background solution in Battefeld:2005cj , as well as the sections on vector and tensor perturbations are unaffected by this revision., e.g. by using the technique introduced in this draft or by including entropy perturbations, as emphasized in Peter:2003 and later on in PintoNeto:2004wf (they derived a regular forth order equation for the full Bardeen potential).
###### Acknowledgements.
We would like to thank the referee of Battefeld:2005cj for drawing our attention to a possible problem associated with the splitting method, and N. Afshordi, D. Battefeld, R. Brandenberger and D. Chung for comments on the draft.
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# Electroweak quasielastic response functions in nuclear matter
## I Introduction
Traditionally, studies of excitations of the nucleus via the interactions between leptons and the nucleus have been centered mostly on the longitudinal ($`R_\text{L}`$) and transverse ($`R_\text{T}`$) nuclear response functions explored in unpolarized, inelastic, inclusive $`(e,e^{})`$ electron scattering. The present paper will also be largely restricted to the investigation of these quantities, although the set of experimentally accessible responses is much larger, including also semi-inclusive responses and responses that arise when initial or final hadronic degrees of freedom are active. In the former, particles are detected in coincidence with the scattered electron $`((e,e^{}p),`$ $`(e,e^{}n),`$ $`(e,e^{}pn),`$ etc.), whereas the latter includes reactions such as $`\stackrel{}{A}(e,e^{}),`$ $`A(e,e^{}\stackrel{}{p}),`$ etc. Of course, in all cases the incident electron may also be polarized. A concerted effort is presently being placed on experimental studies of this extended set of responses where special sensitivities to normally hidden aspects of nuclear structure are expected to exist. Such studies are very challenging and thus only relatively recently have they been made feasible by advances in accelerators, in developments of polarized electron beams and polarized nuclear targets, and in the construction of the required hadron polarimetry. Issues still surround the inclusive unpolarized responses themselves, however, and since a successful level of understanding of the underlying nuclear dynamics would be incomplete without a coherent picture of the entire set of responses, unpolarized inclusive and semi-inclusive/polarized, the former still deserve continued study and so provide the focus for the present article.
Beyond the electromagnetic (EM: parity-conserving, vector) responses $`R_\text{L}`$ and $`R_\text{T}`$ our study will also include their parity-violating analogs $`R_\text{L}^{\text{AV}}`$ and $`R_\text{T}^{\text{AV}}`$ as well as the nuclear parity-violating axial response $`R_\text{T}^{}^{\text{VA}}`$. Here AV indicates that the axial leptonic and vector hadronic currents enter; VA indicates the converse. This larger set of inclusive responses may be explored via the inclusive scattering of longitudinally polarized electrons from unpolarized nuclei, since the electron helicity asymmetry is parity-violating. The three new responses all arise from interferences between the weak neutral current (WNC) and EM current. In the cases of $`R_\text{L}^{\text{AV}}`$ and $`R_\text{T}^{\text{AV}}`$ it is the vector part of the WNC that enters and this is believed to be closely related to the EM current (in the absence of strangeness content in the nucleus these two responses are tied to $`R_\text{L}`$ and $`R_\text{T}`$); in the case of $`R_\text{T}^{}^{\text{VA}}`$ it is the axial part of the WNC that enters, namely an interesting new inclusive nuclear response function.
The vector responses (the four labeled either L or T) are, of course, interesting in their own right, since disentangling them through measurements of both parity-conserving and -violating inclusive electron scattering would permit the isolation of the isoscalar and isovector contributions they contain, as will be discussed in detail later. Accomplishing this separation would represent a significant step forward in our understanding of nuclear structure: indeed, for a long time researchers have sought possible ways of “measuring” how nuclear correlations work in the isoscalar and isovector channels.
A further point worth noting is that analogous responses also play a role in the scattering of hadrons from nuclei. However, in order to interpret the experimental data properly, in addition to the response functions there one also needs a reliable description of the reaction mechanism. In fact, unlike either real or virtual photons, hadrons are mostly absorbed or scattered at the surface of the nucleus. Moreover, the hadrons disrupt the nucleus to a much larger extent than do photons, and therefore the interpretation of reactions induced by photons is generally felt to be under better control than those induced by hadrons. Moderating this statement to some degree is the fact that electron scattering is still somewhat flawed by a not yet fully satisfactory understanding of dispersive effects and of the distortion of the electron waves moving in the EM potential of the nucleus. Accurate knowledge of the nuclear response functions gained with electron scattering provides a way to test the reaction mechanisms of the models employed in hadron scattering. However, the electroweak studies do not provide all of the information we seek and in this regard it is worth noting that in some cases reactions induced by hadrons give access to nuclear responses that are not easily extracted in electron scattering, the best example in this connection being offered by the long sought after spin-longitudinal isovector response.
Turning now to the problem of modeling the nuclear response functions, in the present work we confine our attention to the quasifree region, which is well suited for a microscopic treatment in terms of nucleons and mesons, specifically in terms of the standard field theoretical techniques that we employ. Although the $`\mathrm{\Delta }`$ peak may also be treated in the same framework, it will not be dealt with here to curb the length of this article. Since we are concerned with kinematical regions where relativistic effects are relevant, a good starting point for the development of our approach is given by the Relativistic Fermi Gas (RFG), a covariant model in the sense that its ingredients are the fully relativistic nucleon propagators and EM/WNC vertices. Of course the RFG misses surface and finite-size effects. First of all, these are of secondary importance in obtaining a general understanding of the scattering of electrons in the quasielastic and $`\mathrm{\Delta }`$ peak domains. Secondly, they can be accounted for within the semiclassical approach, which exploits the advantages offered by the translational invariance of the RFG and yet is able to incorporate some of the physics of a finite system.
As discussed in detail later, the perturbative approach we follow requires the setting up of a nuclear mean field (Hartree-Fock, HF) and the treatment of the residual interaction effects in the fully antisymmetrized Random Phase Approximation (RPA). In addition one needs as preliminary input the nucleon-nucleon force: since we would like to view the nucleus as an interacting system of baryons and mesons, a natural choice in this connection is given by a meson-exchange interaction such as the Bonn potential, which can be cast in the framework of an effective field theory. Another preliminary problem relates to the short-range nuclear correlations induced by the violent repulsion present in the nucleon-nucleon force at small distances. A technique for their treatment is indeed available, namely the summation of the Brueckner ladder diagrams; however, it is not yet possible to employ this technique covariantly especially at high density, where on the one hand ladder diagrams are increasingly important and on the other the role of relativity cannot be ignored. In lieu of this, in a few of the results discussed in the following section we shall indicate what insight can be gained by employing a parameterization of a non-relativistic $`G`$-matrix based on the Bonn potential.
## II Quasielastic response functions for inclusive electron scattering
As mentioned in the Introduction, in past years quasielastic electron scattering from nuclei has been the subject of intense experimental and theoretical (see, e. g., Refs. ) investigations. The first aim of the theoretical studies is to test the available nuclear models; once the nuclear physics issues are well understood, one might then hope to gain insight into other aspects of the problem, for instance into the form factors of the nucleon, which can be extracted from the data with an accuracy that is strictly connected to our ability to handle the nuclear physics.
In principle, the quasifree regime is thought to be the obvious place to focus on, as one hopes that there the physical quantities of interest may be computed in a reliable way, while also in this case in practice one has to cope with significant problems. Many diverse techniques have been employed in the literature. Each of them has its own relative merits and deficiencies and clearly it would be highly desirable to be able to reach some degree of convergence in their outcomes.
In the following , we shall be concerned with Green’s function techniques as introduced, e. g., in Ref. . This method can be, and has been, applied both to finite nuclei and nuclear matter. Here, we shall focus on nuclear matter, having in mind applications to electron scattering (that is, without the complications introduced by the reaction mechanism of hadronic probes) in a range from a few hundreds to about 1 GeV/c of transferred momentum where the quasielastic peak is far from low-energy resonances and not too much affected by finite-size effects. The use of nuclear matter reduces the computational load, thus allowing a more straightforward implementation of more sophisticated theoretical schemes than would otherwise be feasible, and this makes it easier to develop and test approximation methods that might subsequently also be utilized for calculations in finite nuclei.
Let us now briefly summarize the theoretical framework that we shall discuss in detail in the following subsections.
A first issue one has to confront in setting up the formalism concerns the treatment of relativistic effects. Kinematical effects, while obviously rather important, can be included in a straightforward way. The treatment of dynamical effects is more delicate. Two main paths have been followed in the literature, either using field theoretical methods (as done, e. g., in the Walecka model and its derivations ) or using potential techniques (i. e., employing phenomenological potentials truncated at some order in the non-relativistic expansion). Here we shall take the second path, but to limit the amount of material to be covered, we shall discuss only non-relativistic potentials.
The extensions necessary to include higher-order relativistic terms are discussed in Ref. , where the nuclear response functions have been calculated using techniques similar to the ones explained below, using as an input the relativistic Bonn potential expanded in powers of $`P/m_N`$ and $`q/m_N`$ up to second order — $`P`$ and $`q`$ being the average of the incoming and outgoing nucleon momenta and the exchanged momentum, respectively. As shown in , the effect of these dynamical relativistic corrections is significant; indeed, the validity of that expansion at high momenta and the inclusion in that framework of short-range nucleon-nucleon correlations has yet to be explored (see, however, Refs. ).
Next, one should choose the phenomenological input potential and, in connection with this choice, attempt to cope with the problem of dealing with short-range correlations. All of the formulae given in the following sections are based on a generic one-boson-exchange potential. They can thus be used both with a bare phenomenological interaction — such as one of the Bonn potential variants — or with a one-boson-exchange parameterization of the $`G`$-matrix generated from some potential. The use of an effective interaction derived from a $`G`$-matrix is a common way of including short-range correlations. However, apart from the relativistic issue, one should be aware of possible problems due to the use of a local potential to fit non-local matrix elements. At least in a few cases discussed in the literature this does not appear to be a reason for concern . On the other hand, possible effects arising only in the quasielastic regime remain completely unexplored. Indeed, $`G`$-matrices employed in quasielastic calculations are usually generated using bound-state boundary conditions, which make them real and practically energy-independent, while in general they are both complex and energy-dependent.
Once we have fixed the effective interaction, we can proceed to consider a hierarchy of approximation schemes.
The lowest-order approximation is, of course, given by the free Fermi gas. Then, one may include mean-field correlations at the HF level (or Brueckner-Hartree-Fock (BHF) if short-range correlations are accounted for). In nuclear matter a HF calculation can be done exactly without too much effort. Later we show how a quite accurate analytic approximation can be derived, and how this is needed to combine the HF and RPA schemes. The latter is the last resummation technique we shall discuss. It should be noticed that even in nuclear matter the calculation of the antisymmetrized RPA response functions is not trivial. Indeed, most calculations, labeled “RPA” in the literature, are actually performed in the so-called “ring approximation”, where only the direct contributions are kept. For this case, in nuclear matter one gets a simple algebraic equation for the response. Here, we use the continued fraction (CF) technique to provide a semi-analytical estimate of the full RPA response (see Refs. and for alternative methods). Calculations with this method have been performed both in finite nuclei and in nuclear matter , always truncating the CF expansion at first order because of the difficulty of the numerical calculations involved. We have pushed the analytical calculation far enough to yield not only a fast and accurate estimate of the first-order CF contribution, but also of the second-order one. Since the rate of convergence of the CF expansion cannot be assessed on the basis of general theorems, this is the only way of getting a quantitative grip on the quality of the approximation. As mentioned before, HF (and kinematical relativistic) effects can then be incorporated in the RPA calculation, yielding as the final approximation scheme a HF-RPA (or BHF-RPA) response function.
Of course, several many-body contributions have been left out in our analysis. However the classes of many-body diagrams discussed here already allow one to capture the main features of the quasielastic response and, since semi-analytical methods have been developed for their computation, our formalism constitutes a valid starting point for the study of other many-body effects.
### A Response functions
We consider an infinite system of interacting nucleons at some density fixed by the Fermi momentum $`k_F`$. For the kinetic energies of the nucleons we can choose either relativistic or non-relativistic expressions, whereas we assume that the interactions take place through a non-relativistic potential. For the latter the following expression in momentum space is assumed
$`V(𝒌)`$ $`=`$ $`V_0(k)+V_\tau (k)𝝉_1𝝉_2+V_\sigma (k)𝝈_1𝝈_2+V_{\sigma \tau }(k)𝝈_1𝝈_2𝝉_1𝝉_2`$ (2)
$`+V_t(k)S_{12}(\widehat{𝒌})+V_{t\tau }(k)S_{12}(\widehat{𝒌})𝝉_1𝝉_2,`$
where $`S_{12}`$ is the standard tensor operator and $`V_\alpha (k)`$ represents the momentum space potential in channel $`\alpha `$. Here $`V_\alpha (k)`$ has the general form of a static one-boson-exchange potential so that in each spin-isospin channel, namely $`(0,\tau ,\sigma ,\sigma \tau ,t,t\tau )`$, it is represented as a sum of contributions from different mesons, $`V_\alpha _iV_\alpha ^{(i)}`$. In the central channels ($`0`$, $`\tau `$, $`\sigma `$, $`\sigma \tau `$) the contribution from any meson can be expressed as the combination of a short-range (“$`\delta `$”) piece and a longer range (“momentum-dependent”) piece<sup>*</sup><sup>*</sup>* The nomenclature stems from the fact that, in the absence of form factors, $`V_\delta `$ is a constant and is represented by a Dirac $`\delta `$-function in coordinate space, whereas $`V_{\text{MD}}`$ is the momentum-dependent piece.:
$`V_\delta ^{(i)}(k)`$ $`=`$ $`g_\delta ^{(i)}\left({\displaystyle \frac{\mathrm{\Lambda }_i^2m_i^2}{\mathrm{\Lambda }_i^2+k^2}}\right)^{\mathrm{}}`$ (4)
$`V_{\text{MD}}^{(i)}(k)`$ $`=`$ $`g_{\text{MD}}^{(i)}{\displaystyle \frac{m_i^2}{m_i^2+k^2}}\left({\displaystyle \frac{\mathrm{\Lambda }_i^2m_i^2}{\mathrm{\Lambda }_i^2+k^2}}\right)^{\mathrm{}},\mathrm{}=0,1,2,`$ (5)
whereas in the tensor channels ($`t`$, $`t\tau `$) one has
$$V_{\text{TN}}^{(i)}(k)=g_{\text{TN}}^{(i)}\frac{k^2}{m_i^2+k^2}\left(\frac{\mathrm{\Lambda }_i^2m_i^2}{\mathrm{\Lambda }_i^2+k^2}\right)^{\mathrm{}},\mathrm{}=0,1,2.$$
(6)
In Eqs. (*II A), $`g_\delta ^{(i)}`$, $`g_{\text{MD}}^{(i)}`$ and $`g_{\text{TN}}^{(i)}`$ are the (dimensional) coupling constants of the $`i`$-th meson, $`m_i`$ is its mass and $`\mathrm{\Lambda }_i`$ the cut-off; more generally, potentials without form factors or with monopole or dipole form factors are allowed.
Our starting point is the Galitskii-Migdal integral equation for the particle-hole (ph) four-point Green’s function Capital letters refer to four-vectors and lower-case letters to three-vectors; the Greek letters $`\alpha ,\beta ,\mathrm{}`$ refer to a set of spin-isospin quantum numbers.,
$`G_{\alpha \beta ,\gamma \delta }^{\text{ph}}(K+Q,K;P+Q,P)=G_{\alpha \gamma }(P+Q)G_{\delta \beta }(P)(2\pi )^4\delta (KP)`$ (7)
$`+iG_{\alpha \lambda }(K+Q)G_{\lambda ^{}\beta }(K){\displaystyle \frac{d^4T}{(2\pi )^4}\mathrm{\Gamma }_{\lambda \lambda ^{},\mu \mu ^{}}^{13}(K+Q,K;T+Q,T)G_{\mu \mu ^{},\gamma \delta }^{\text{ph}}(T+Q,T;P+Q,P)},`$ (8)
(9)
diagrammatically illustrated in Fig. 1. In Eq. (9), $`G`$ represents the exact one-body Green’s function, whereas $`\mathrm{\Gamma }^{13}`$ is the irreducible vertex function in the ph channel.
Given $`G^{\text{ph}}`$ one can then define the polarization propagator
$`\mathrm{\Pi }_{\alpha \beta ,\gamma \delta }(Q)`$ $``$ $`\mathrm{\Pi }_{\alpha \beta ,\gamma \delta }(q,\omega )`$ (10)
$`=`$ $`i{\displaystyle \frac{d^4P}{(2\pi )^4}\frac{d^4K}{(2\pi )^4}G_{\alpha \beta ,\gamma \delta }^{\text{ph}}(K+Q,K;P+Q,P)},`$ (11)
whose diagrammatic representation is displayed in Fig. 2. Note that for $`\mathrm{\Pi }(q,\omega )`$ one cannot in general write down an integral (or algebraic) equation.
In the case of electron scattering, one can define charge — or longitudinal — and magnetic — or transverse — polarization propagators. These, in the non-relativistic regime, read
$`\mathrm{\Pi }_\text{L}^I(q,\omega )=\text{tr}[\widehat{O}_\text{L}^I\widehat{\mathrm{\Pi }}(q,\omega )\widehat{O}_\text{L}^I]`$ (13)
$`\mathrm{\Pi }_\text{T}^I(q,\omega )={\displaystyle \underset{ij}{}}\mathrm{\Lambda }_{ji}\mathrm{\Pi }_{ij}^I(q,\omega ),\mathrm{\Pi }_{ij}^I(q,\omega )=\text{tr}[\widehat{O}_{\text{T};i}^I\widehat{\mathrm{\Pi }}(q,\omega )\widehat{O}_{\text{T};j}^I]`$ (14)
$`\mathrm{\Lambda }_{ij}=(\delta _{ij}𝒒_i𝒒_j/q^2)/2,`$ (15)
where, for brevity, the dependence upon the spin-isospin indices has been represented in matrix form, introducing hats to indicate matrices. In Eqs. (2), $`I`$ labels the isospin channel and the longitudinal and transverse vertex operators are given by:
$$\{\begin{array}{c}\widehat{O}_\text{L}^{I=0}=1/2\hfill \\ \widehat{O}_\text{L}^{I=1}=\tau _3/2\hfill \end{array}\{\begin{array}{c}\widehat{O}_{\text{T};i}^{I=0}=\sigma _i/2\hfill \\ \widehat{O}_{\text{T};i}^{I=1}=\sigma _i\tau _3/2.\hfill \end{array}$$
(16)
The inelastic inclusive scattering cross section where the momentum $`q`$ and energy $`\omega `$ are transferred to the nucleus is a linear combination of the imaginary parts of $`\mathrm{\Pi }_{\text{L,T}}(q,\omega )`$. It is then customary to define longitudinal and transverse response functions according to
$$R_{\text{L,T}}(q,\omega )=R_{\text{L,T}}^{I=0}(q,\omega )+R_{\text{L,T}}^{I=1}(q,\omega ),$$
(17)
which are related to $`\mathrm{\Pi }_{\text{L,T}}`$ by
$`R_{\text{L,T}}^I(q,\omega )`$ $`=`$ $`{\displaystyle \frac{V}{\pi }}f_{\text{L,T}}^{(I)}{}_{}{}^{2}(q,\omega )\text{Im}\mathrm{\Pi }_{\text{L,T}}^I(q,\omega )`$ (18)
$`=`$ $`{\displaystyle \frac{3\pi A}{2k_F^3}}f_{\text{L,T}}^{(I)}{}_{}{}^{2}(q,\omega )\text{Im}\mathrm{\Pi }_{\text{L,T}}^I(q,\omega ),`$ (19)
where $`V`$ is the volume, $`A`$ the mass number and the $`f_{\text{L,T}}^{(I)}{}_{}{}^{2}`$ embody the squared EM form factors of the nucleon. The latter are briefly discussed in Appendix A.
### B Non-relativistic vs relativistic kinematics
The response functions introduced above have been defined as functions of the momentum transfer $`q`$ and energy transfer $`\omega `$. Actually, it is possible — and convenient — to define a scaling variable $`\psi `$ that is a function of $`q`$ and $`\omega `$ and may be used in place of $`\omega `$. This variable is such that the responses of a free Fermi gas in the non-Pauli-blocked region ($`q>2k_F`$) can be expressed in terms of the variable $`\psi `$ only (apart from $`q`$-dependent multiplicative factors). We shall see that even in the Pauli-blocked region and for an interacting system it is convenient to use the pair of variables ($`q`$,$`\psi `$) instead of ($`q`$,$`\omega `$).
Besides the obvious advantages related to the use of a scaling variable (see Section IV), there is another reason for expressing the responses in terms of $`\psi `$: when the latter is used the responses viewed as functions of $`\psi `$ turn out to adjust to the form assumed for the nucleon kinetic energy. To be more specific, starting from either a non-relativistic or relativistic Fermi gas, one is always led to essentially the same dependence of the responses upon the corresponding $`\psi `$ variable.
We shall see in the following subsections that the energy denominators of the free nucleon propagators appearing in the Feynman diagrams for the response functions are always given by $`\omega ϵ_{𝒌+𝒒}^{(0)}+ϵ_𝒌^{(0)}`$, where $`ϵ_𝒌^{(0)}`$ is the kinetic energy of a nucleon of momentum $`k`$ and $`k<k_F`$. In the non-relativistic case
$`\omega ϵ_{𝒌+𝒒}^{(0)\text{nr}}+ϵ_𝒌^{(0)\text{nr}}`$ $`=`$ $`\omega {\displaystyle \frac{(𝒌+𝒒)^2}{2m_N}}+{\displaystyle \frac{k^2}{2m_N}}`$ (20)
$`=`$ $`{\displaystyle \frac{qk_F}{m_N}}\left(\psi _{\text{nr}}\widehat{𝒒}{\displaystyle \frac{𝒌}{k_F}}\right),`$ (21)
where
$$\psi _{\text{nr}}=\frac{1}{k_F}\left(\frac{\omega m_N}{q}\frac{q}{2}\right)$$
(22)
is the standard scaling variable of the non-relativistic Fermi gas and $`m_N`$ the nucleon mass.
In the relativistic case, one would have
$$\omega ϵ_{𝒌+𝒒}^{(0)\text{r}}+ϵ_𝒌^{(0)\text{r}}=\omega \sqrt{(𝒌+𝒒)^2+m_N^2}+\sqrt{k^2+m_N^2};$$
(23)
however, in Ref. it was shown that at the pole (where the above vanishes) a very good approximation for Eq. (23) obtains by using Eq. (21) with, instead of $`\psi _{\text{nr}}`$,
$$\psi _\text{r}=\frac{1}{k_F}\left[\frac{\omega m_N(1+\omega /2m_N)}{q}\frac{q}{2}\right]$$
(24)
and then by multiplying the free response by $`1+\omega /m_N`$, which is proportional to the Jacobian of the transformation from the variable $`\omega `$ to the variable $`\psi `$. Thus the use of the scaling variable in Eq. (24) entails the substitution
$$\omega ϵ_{𝒌+𝒒}^{(0)\text{r}}+ϵ_𝒌^{(0)\text{r}}\omega \left(1+\frac{\omega }{2m_N}\right)ϵ_{𝒌+𝒒}^{(0)\text{nr}}+ϵ_𝒌^{(0)\text{nr}}.$$
(25)
In turn, this implies that the pole (which provides the contribution to the imaginary part of the propagator) is located at $`\omega =\sqrt{m_N^2+q^2+2𝒒𝒌}m_N`$, namely at the place predicted by the exact expression in Eq. (23) when $`k^2`$ is neglected with respect to $`m_N^2`$. As stated above, since $`k`$ is always below $`k_F`$, this is a good approximation and, indeed, the free RFG response calculated using the scaling variable in Eq. (24) reproduces that of the exact calculation accurately, the discrepancy being typically below 1%.
However, in the calculation of higher-order (RPA) contributions, the real part of the energy denominators also comes into play and the validity of the approximation far from the pole should also be checked. With some algebra — and assuming $`k^2/m_N^21`$ — one can write
$`\omega ϵ_{𝒌+𝒒}^{(0)\text{r}}+ϵ_𝒌^{(0)\text{r}}`$ $``$ $`{\displaystyle \frac{qk_F}{m_N}}{\displaystyle \frac{\psi _r\widehat{𝒒}𝒌/k_F}{{\displaystyle \frac{1}{2}}\left(1+{\displaystyle \frac{\omega }{m_N}}+\sqrt{1+{\displaystyle \frac{q^2+2𝒒𝒌}{m_N^2}}}\right)}}`$ (26)
$``$ $`{\displaystyle \frac{qk_F}{m_N}}{\displaystyle \frac{\psi _r\widehat{𝒒}𝒌/k_F}{1+\omega /m_N}},`$ (27)
where, in the last passage, we have replaced the square root with its value at the pole. In Fig. 3, we display the real part of the free polarization propagator (defined in the following subsection) using the exact relativistic dispersion relation and the prescription of Eq. (27) at $`q=500`$ MeV/c and 1 GeV/c as a function of $`\omega `$. The agreement between the two ways of calculating $`\text{Re}\mathrm{\Pi }^{(0)}`$ is quite good at both momenta.
Equation (27) provides an approximation for the free ph propagator. A prescription to obtain the (kinematically) relativistic polarization propagators at any order in the RPA expansion (see Section II E) can easily be obtained by noting that $`\mathrm{\Pi }^{(n)}`$ — the $`n`$-th order contribution to the RPA chain — contains $`n+1`$ ph propagators; one then has
$$\mathrm{\Pi }^{(n)\text{r}}(q,\omega )=\left(1+\frac{\omega }{m_N}\right)^{n+1}\mathrm{\Pi }^{(n)\text{nr}}(q,\omega (1+\omega /2m_N)).$$
(29)
Actually, all of the response functions derived below are expressed in terms of a generic scaling variable $`\psi `$, as $`\mathrm{\Pi }^{(n)}(q,\psi )`$. One can then get the non-relativistic response by using the (exact) expression in Eq. (22) for $`\psi `$ and the relativistic response by using the (approximate) form in Eq. (24) and multiplying each polarization propagator by the appropriate power of $`1+\omega /m_N`$, i. e.
$$\mathrm{\Pi }^{(n)\text{r}}(q,\omega )=\left(1+\frac{\omega }{m_N}\right)^{n+1}\mathrm{\Pi }^{(n)\text{nr}}(q,\psi _\text{r}).$$
(30)
Note that in the calculations of Refs. only an overall Jacobian factor, $`1+\omega /m_N`$, has been applied to the RPA response functions. In typical kinematical conditions the size of the error introduced by this further approximation is of the order of a few percent.
### C Free response
Although the free Fermi gas response function is a subject for textbooks (see, e. g., Ref. ), it is useful to derive it here using a slightly different approach, since it illustrates at the simplest level the method we have adopted to overcome a technical difficulty one meets in nuclear matter calculations — namely the presence of $`\theta `$-functions, which considerably complicates analytic integrations. As a side effect, the expression for $`\mathrm{\Pi }^{(0)}`$ also comes out to be much more compact than in standard treatments.
From Eqs. (2) and (16), one immediately finds that
$$\mathrm{\Pi }_{\text{L};I=0}^{(0)}=\mathrm{\Pi }_{\text{L};I=1}^{(0)}=\mathrm{\Pi }_{\text{T};I=0}^{(0)}=\mathrm{\Pi }_{\text{T};I=1}^{(0)}\mathrm{\Pi }^{(0)},$$
(31)
where following Eqs. (9) and (11) we have defined
$$\mathrm{\Pi }^{(0)}(q,\omega )=\frac{d𝒌}{(2\pi )^3}G_{\text{ph}}^{(0)}(𝒌,𝒒;\omega ),$$
(32)
having set
$$G_{\text{ph}}^{(0)}(𝒌,𝒒;\omega )=i\frac{dk_0}{2\pi }G^{(0)}(𝒌+𝒒,k_0+\omega )G^{(0)}(k,k_0),$$
(33)
$`G^{(0)}(k,k_0)`$ being the free one-body propagator
$$G^{(0)}(k,k_0)=\frac{\theta (kk_F)}{k_0ϵ_𝒌^{(0)}+i\eta }+\frac{\theta (k_Fk)}{k_0ϵ_𝒌^{(0)}i\eta }.$$
(34)
The integration over $`k_0`$ in Eq. (33) is straightforward, yielding
$$G_{\text{ph}}^{(0)}(𝒌,𝒒;\omega )=\frac{\theta (k_Fk)\theta (|𝒌+𝒒|k_F)}{\omega ϵ_{𝒌+𝒒}^{(0)}+ϵ_𝒌^{(0)}+i\eta }+\frac{\theta (kk_F)\theta (k_F|𝒌+𝒒|)}{\omega +ϵ_{𝒌+𝒒}^{(0)}ϵ_𝒌^{(0)}+i\eta },$$
(35)
which, inserted back into Eq. (32), would give the standard definition of $`\mathrm{\Pi }^{(0)}`$. Instead, let us rewrite $`G_{\text{ph}}^{(0)}`$ as
$`G_{\text{ph}}^{(0)}(𝒌,𝒒;\omega )`$ $`=`$ $`{\displaystyle \frac{\theta (k_Fk)\theta (|𝒌+𝒒|k_F)}{\omega ϵ_{𝒌+𝒒}^{(0)}+ϵ_𝒌^{(0)}+i\eta }}+{\displaystyle \frac{\theta (kk_F)\theta (k_F|𝒌+𝒒|)}{\omega +ϵ_{𝒌+𝒒}^{(0)}ϵ_𝒌^{(0)}+i\eta }}`$ (37)
$`+{\displaystyle \frac{\theta (k_Fk)\theta (k_F|𝒌+𝒒|)}{\omega ϵ_{𝒌+𝒒}^{(0)}+ϵ_𝒌^{(0)}+i\eta _\omega }}+{\displaystyle \frac{\theta (k_Fk)\theta (k_F|𝒌+𝒒|)}{\omega +ϵ_{𝒌+𝒒}^{(0)}ϵ_𝒌^{(0)}i\eta _\omega }},`$
having added and subtracted the quantity in the second line, where we have set $`\eta _\omega =\text{sign}(\omega )\eta `$. A few algebraic manipulations then yield
$$G_{\text{ph}}^{(0)}(𝒌,𝒒;\omega )=\frac{\theta (k_Fk)\theta (k_F|𝒌+𝒒|)}{\omega ϵ_{𝒌+𝒒}^{(0)}+ϵ_𝒌^{(0)}+i\eta _\omega }.$$
(38)
Hence, from Eq. (32) one gets
$`\mathrm{\Pi }^{(0)}(q,\omega )`$ $`=`$ $`{\displaystyle \frac{d𝒌}{(2\pi )^3}\theta (k_Fk)\left[\frac{1}{\omega ϵ_{𝒌+𝒒}^{(0)}+ϵ_𝒌^{(0)}+i\eta _\omega }+\frac{1}{\omega ϵ_{𝒌+𝒒}^{(0)}+ϵ_𝒌^{(0)}i\eta _\omega }\right]}`$ (39)
$`=`$ $`{\displaystyle \frac{m_N}{q}}{\displaystyle \frac{k_F^2}{(2\pi )^2}}\left[𝒬^{(0)}(\psi )𝒬^{(0)}(\psi +\overline{q})\right].`$ (40)
Note that only one $`\theta `$-function forcing $`k`$ below $`k_F`$ is left, Pauli blocking being enforced by cancellations between the energy denominators. In Eq. (40), we have introduced $`\overline{q}=q/k_F`$ and the dimensionless function
$$𝒬^{(0)}(\psi )=\frac{1}{2}_1^1𝑑y\frac{1y^2}{\psi y+i\eta _\omega },$$
(41)
which is easily evaluated, yielding
$`\text{Re}𝒬^{(0)}(\psi )`$ $`=`$ $`\psi +{\displaystyle \frac{1}{2}}(1\psi ^2)\text{ln}\left|{\displaystyle \frac{1+\psi }{1\psi }}\right|={\displaystyle \frac{2}{3}}[Q_0(\psi )Q_2(\psi )]`$ (43)
$`\text{Im}𝒬^{(0)}(\psi )`$ $`=`$ $`\text{sign}(\omega )\theta (1\psi ^2){\displaystyle \frac{\pi }{2}}(1\psi ^2)=\text{sign}(\omega )\theta (1\psi ^2){\displaystyle \frac{\pi }{3}}[P_0(\psi )P_2(\psi )],`$ (44)
where $`P_n`$ and $`Q_n`$ are Legendre polynomials and Legendre functions of the second kind, respectively.
The expression in Eq. (40) has a simple physical interpretation. If one switches off Pauli blocking, the response of a Fermi sphere, with four particles per momentum state up to $`k_F`$, is given by a parabola over the response region $`q^2/2m_Nqk_F/m_N<\omega <q^2/2m_N+qk_F/m_N`$, that is, the curve obtained joining the dotted line and the parabolic section of the solid line in Fig. 4. With respect to the Pauli blocking, two kinds of spurious terms arise when $`k`$ and $`|𝒌+𝒒|`$ are both below the Fermi surface. If $`|𝒌+𝒒|>k`$, then a spurious contribution occurs in the Pauli-forbidden region $`0<\omega <qk_F/m_Nq^2/2m_N`$, whereas if $`|𝒌+𝒒|<k`$, then a contribution occurs with the same strength for $`q^2/2m_Nqk_F/m_N<\omega <0`$. Hence, in order to get the correct response function, one simply subtracts — for a given $`\omega >0`$ in the Pauli-forbidden region — the total of the spurious contributions at $`\omega `$, thus getting the familiar linear dependence on $`\omega `$. Graphically, as illustrated in Fig. 4, this amounts to reflecting the response at negative transferred energies in the vertical axis and then subtracting it.
### D Hartree-Fock response
The HF polarization propagator in nuclear matter is obtained by dressing the one-body propagators appearing in $`\mathrm{\Pi }^{(0)}`$ with the first-order self-energy $`\mathrm{\Sigma }^{(1)}`$, so that one can follow essentially the same derivation of the previous subsection. The spin-isospin matrix elements are the same as for the free response, yielding
$$\mathrm{\Pi }_{\text{L};I=0}^{\text{HF}}=\mathrm{\Pi }_{\text{L};I=1}^{\text{HF}}=\mathrm{\Pi }_{\text{T};I=0}^{\text{HF}}=\mathrm{\Pi }_{\text{T};I=1}^{\text{HF}}\mathrm{\Pi }^{\text{HF}},$$
(45)
where
$$\mathrm{\Pi }^{\text{HF}}(q,\omega )=\frac{d𝒌}{(2\pi )^3}G_{\text{ph}}^{\text{HF}}(𝒌,𝒒;\omega )$$
(46)
and
$$G_{\text{ph}}^{\text{HF}}(𝒌,𝒒;\omega )=i_{\mathrm{}}^{\mathrm{}}\frac{dk_0}{2\pi }G^{\text{HF}}(𝒌+𝒒,k_0+\omega )G^{\text{HF}}(k,k_0),$$
(47)
$`G^{\text{HF}}(k,k_0)`$ being the HF one-body propagator
$$G^{\text{HF}}(k,k_0)=\frac{\theta (kk_F)}{k_0ϵ_𝒌^{(1)}+i\eta }+\frac{\theta (k_Fk)}{k_0ϵ_𝒌^{(1)}i\eta },$$
(48)
with
$$ϵ_𝒌^{(1)}=ϵ_𝒌^{(0)}+\mathrm{\Sigma }^{(1)}(k).$$
(49)
Since the first-order self-energy does not depend on the energy, the integration over $`k_0`$ can be performed along the lines of Eqs. (35)–(38), yielding
$$G_{\text{ph}}^{\text{HF}}(𝒌,𝒒;\omega )=\frac{\theta (k_Fk)\theta (k_F|𝒌+𝒒|)}{\omega ϵ_{𝒌+𝒒}^{(1)}+ϵ_𝒌^{(1)}+i\eta _\omega }$$
(50)
and, finally,
$$\mathrm{\Pi }^{\text{HF}}(q,\omega )=\frac{d𝒌}{(2\pi )^3}\theta (k_Fk)\left[\frac{1}{\omega ϵ_{𝒌+𝒒}^{(1)}+ϵ_𝒌^{(1)}+i\eta _\omega }+\frac{1}{\omega ϵ_{𝒌+𝒒}^{(1)}+ϵ_𝒌^{(1)}i\eta _\omega }\right].$$
(51)
The HF response function is proportional to the imaginary part of $`\mathrm{\Pi }^{\text{HF}}`$:
$`\text{Im}\mathrm{\Pi }^{\text{HF}}(q,\omega )`$ $`=`$ $`\text{sign}(\omega )\pi {\displaystyle \frac{d𝒌}{(2\pi )^3}\theta (k_Fk)\left[\delta \left(\omega ϵ_{𝒌+𝒒}^{(1)}+ϵ_𝒌^{(1)}\right)\delta \left(\omega ϵ_{𝒌+𝒒}^{(1)}+ϵ_𝒌^{(1)}\right)\right]}`$ (52)
$`=`$ $`\text{sign}(\omega )\pi {\displaystyle \frac{m_N}{q}}{\displaystyle \frac{1}{(2\pi )^2}}`$ (54)
$`\times {\displaystyle _0^{k_F}}dkk{\displaystyle \frac{1}{m_N}}[m_N^{}(\sqrt{k^2+q^2+2qy_0})m_N^{}(\sqrt{k^2+q^2+2q\overline{y}_0})],`$
having defined the effective mass as
$$m_N^{\text{nr}}(k)=\frac{m_N}{1+{\displaystyle \frac{m_N}{k}}{\displaystyle \frac{d\mathrm{\Sigma }^{(1)}}{dk}}}$$
(56)
or
$$m_N^\text{r}(k)=\frac{\sqrt{m_N^2+k^2}}{1+{\displaystyle \frac{\sqrt{m_N^2+k^2}}{k}}{\displaystyle \frac{d\mathrm{\Sigma }^{(1)}}{dk}}},$$
(57)
for the non-relativistic or relativistic case, respectively, whereas $`y_0`$ and $`\overline{y}_0`$ solve the equations
$$\{\begin{array}{ccc}f_{\text{HF}}(\omega |k,y_0)& =& 0\\ f_{\text{HF}}(\omega |k,\overline{y}_0)& =& 0,\end{array}$$
(58)
with
$`f_{\text{HF}}^{\text{nr}}(\omega |k,y)`$ $`=`$ $`\omega {\displaystyle \frac{q^2}{2m_N}}{\displaystyle \frac{qy}{m_N}}\mathrm{\Sigma }^{(1)}(\sqrt{k^2+q^2+2qy})+\mathrm{\Sigma }^{(1)}(k)`$ (60)
$`f_{\text{HF}}^\text{r}(\omega |k,y)`$ $`=`$ $`\omega \sqrt{m_N^2+k^2+q^2+2qy}+\sqrt{m_N^2+k^2}{\displaystyle \frac{qy}{m_N}}`$ (62)
$`\mathrm{\Sigma }^{(1)}(\sqrt{k^2+q^2+2qy})+\mathrm{\Sigma }^{(1)}(k).`$
Although the evaluation of the HF response is numerically quite straightforward, in Ref. an analytic approximation for $`\text{Im}\mathrm{\Pi }^{\text{HF}}`$ has been worked out, with the aim of using it to include the HF field in RPA calculations. Here, it will be shown that the analytic approximation is valid not only for the HF response, but more generally, although in the HF case one can directly assess the good accuracy of the procedure.
In any Feynman diagram considered here and in the following, the nucleon self-energy enters through the ph energy denominators,
$$\omega ϵ_{𝒌+𝒒}^{(1)\text{nr}}+ϵ_𝒌^{(1)\text{nr}}=\omega \frac{(𝒌+𝒒)^2}{2m_N}+\frac{k^2}{2m_N}\mathrm{\Sigma }^{(1)}(|𝒌+𝒒|)+\mathrm{\Sigma }^{(1)}(k),$$
(63)
where the non-relativistic expression for the nucleon kinetic energy has been used. In Eq. (63), one can always assume that $`k<k_F`$ and $`|𝒌+𝒒|>k_F`$. Although the latter inequality is not immediately apparent from, e. g., Eq. (51), remember that cancellations between the energy denominators are such as to enforce the Pauli principle; the same will also be true for the RPA diagrams It should also be noted that the infinite Fermi gas is more in touch with the physics for relatively large momenta ($`q2k_F`$), where the above conditions are satisfied by definition..
Clearly, if $`\mathrm{\Sigma }^{(1)}(k)`$ were parabolic in the momentum, the inclusion of the self-energy would be achieved simply by substituting an effective mass for $`m_N`$. For realistic potentials, a parabolic fit for the self-energy over the whole range of momenta is in general not a good approximation. It is a good approximation, on the other hand, to fit separately the particle and hole parts of the self-energy, the fit being restricted to the range of momenta actually involved in the integration. Since in Eq. (51) (but also in the RPA diagrams discussed later) $`k`$ is integrated from 0 to $`k_F`$ and, furthermore, $`|𝒌+𝒒|>k_F`$, one can set
$`\mathrm{\Sigma }^{(1)}`$ $``$ $`\overline{A}+\overline{B}{\displaystyle \frac{k^2}{2m_N}},0<k<k_F,`$ (64)
$`\mathrm{\Sigma }^{(1)}`$ $``$ $`A+B{\displaystyle \frac{k^2}{2m_N}},\text{max}(qk_F,k_F)<k<q+k_F.`$ (66)
Inserting this “biparabolic approximation” back into Eq. (63), and setting $`\epsilon =\overline{A}A`$ and $`m_N^{\text{nr}}=m_N/(1+B)`$, one gets
$`\omega ϵ_{𝒌+𝒒}^{(1)\text{nr}}+ϵ_𝒌^{(1)\text{nr}}`$ $``$ $`\omega (1+B){\displaystyle \frac{q^2}{2m_N}}(1+B){\displaystyle \frac{𝒒𝒌}{m_N}}+\overline{A}A+(\overline{B}B){\displaystyle \frac{k^2}{2m_N}}`$ (67)
$`=`$ $`{\displaystyle \frac{qk_F}{m_N^{\text{nr}}}}\left\{{\displaystyle \frac{1}{k_F}}\left[(\omega +\epsilon ){\displaystyle \frac{m_N^{\text{nr}}}{q}}{\displaystyle \frac{q}{2}}\right]\widehat{𝒒}{\displaystyle \frac{𝒌}{k_F}}+{\displaystyle \frac{\overline{B}B}{1+B}}{\displaystyle \frac{k_F}{2q}}\left({\displaystyle \frac{k}{k_F}}\right)^2\right\}`$ (68)
$``$ $`{\displaystyle \frac{qk_F}{m_N^{\text{nr}}}}\left[\psi _{\text{nr}}^{}\widehat{𝒒}{\displaystyle \frac{𝒌}{k_F}}\right].`$ (69)
To go from the second to the last line in Eq. (69), we have neglected the term proportional to $`k^2`$, which is expected to be small, since $`k<k_F`$ and, typically, $`q>k_F`$. However, this approximation depends upon the interaction and one should check its validity, since it affects both the parameters $`B`$ and $`\overline{B}`$. In Ref. the term neglected has been shown to be small for the Bonn potential; the same turns out to be true also for the effective interaction employed in the next section.
Equation (69) is similar to the expression (21) for the free energy denominator, but for the substitutions
$`m_N`$ $``$ $`m_N^{\text{nr}}={\displaystyle \frac{m_N}{1+B}}`$ (70)
$`\psi _{\text{nr}}`$ $``$ $`\psi _{\text{nr}}^{}={\displaystyle \frac{1}{k_F}}\left[(\omega +\epsilon ){\displaystyle \frac{m_N^{}}{q}}{\displaystyle \frac{q}{2}}\right]`$ (72)
$`={\displaystyle \frac{\psi _{\text{nr}}+\chi }{1+B}},\chi ={\displaystyle \frac{1}{k_F}}\left({\displaystyle \frac{\epsilon m_N}{q}}B{\displaystyle \frac{q}{2}}\right),`$
(or $`\omega \omega +\epsilon `$).
In Ref. relativistic kinematics had been accounted for by applying to the above formulae the substitution $`\omega \omega (1+\omega /2m_N)`$ previously discussed. The correct approximation can be worked out by starting again from the ph propagator by defining ($`\mathrm{\Delta }\mathrm{\Sigma }^{(1)}(𝒌,𝒒)\mathrm{\Sigma }^{(1)}(k)\mathrm{\Sigma }^{(1)}(|𝒌+𝒒|)`$) and rewriting it as
$`{\displaystyle \frac{1}{\omega ϵ_{𝒌+𝒒}^{(1)\text{r}}+ϵ_𝒌^{(1)\text{r}}}}=`$ (73)
$`={\displaystyle \frac{\omega +\sqrt{k^2+m_N^2}+\mathrm{\Delta }\mathrm{\Sigma }^{(1)}(𝒌,𝒒)+\sqrt{(𝒌+𝒒)^2+m_N^2}}{\omega ^2+2\omega \sqrt{k^2+m_N^2}+2(\omega +\sqrt{k^2+m_N^2})\mathrm{\Delta }\mathrm{\Sigma }^{(1)}(𝒌,𝒒)+[\mathrm{\Delta }\mathrm{\Sigma }^{(1)}(𝒌,𝒒)]^2q^22𝒒𝒌}}`$ (74)
$`{\displaystyle \frac{m_N^\text{r}}{qk_F}}{\displaystyle \frac{1+\omega /m_N+\mathrm{\Delta }^{(1)}/m_N}{\psi _\text{r}^{}\widehat{𝒒}𝒌/k_F}},`$ (75)
where
$`m_N^\text{r}`$ $`=`$ $`{\displaystyle \frac{m_N}{1+B(1+\omega /m_N+\mathrm{\Delta }^{(1)}/m_N)}}`$ (76)
$`\psi _\text{r}^{}`$ $`=`$ $`{\displaystyle \frac{\psi _\text{r}+\chi [1+B(1+\omega /m_N+\mathrm{\Delta }^{(1)}/2m_N)]}{1+B(1+\omega /m_N+\mathrm{\Delta }^{(1)}/m_N)}}`$ (77)
$`\mathrm{\Delta }^{(1)}`$ $`=`$ $`\epsilon B{\displaystyle \frac{q^2}{2m_N}}{\displaystyle \frac{qk_F}{m_N}}\chi ,`$ (78)
with $`\chi `$ already defined in Eq. (72). In deriving Eq. (75), we have assumed that $`k^2m_N^2`$, have evaluated the numerator at the pole thus discarding any angular dependence and, in the denominator, have retained only terms at most linear in $`\widehat{𝒒}𝒌/k_F`$. As one can see, besides the transformation $`\omega \omega (1+\omega /2m_N)`$ there are other relativistic corrections, both to the effective scaling variable and to the Jacobian.
The quality of the approximations introduced above is good: indeed the HF response is reproduced with at most a few percent discrepancy (except on the borders of the response region, where the Fermi gas is anyway unrealistic). Thus, we see that in either the non-relativistic or relativistic case, the prescription to include HF correlations in a response function is simply to replace $`\psi `$ with $`\psi ^{}`$ and $`m_N`$ with $`m_N^{}`$ (and to multiply by a normalization factor when employing relativistic kinematics (see Eqs. (3)). For instance, from Eq. (40) one gets
$$\mathrm{\Pi }^{\text{HF}}(q,\omega )J\frac{m_N^{}}{q}\frac{k_F^2}{(2\pi )^2}\left[𝒬^{(0)}(\psi ^{})𝒬^{(0)}(\psi ^{}+\overline{q})\right],$$
(79)
with
$$J_{\text{nr}}=1$$
(80)
and
$$J_\text{r}=1+\frac{\omega }{m_N}+\frac{\mathrm{\Delta }^{(1)}}{m_N}.$$
(81)
In Appendix B, we give the explicit expressions for the first-order self-energy, based on the generic potential in Eqs. (2)–(*II A).
### E Random phase approximation response
If in Eq. (9) one substitutes the irreducible vertex function $`\mathrm{\Gamma }^{13}`$ with the matrix elements of the bare potential, one gets the so-called random phase approximation to $`G^{\text{ph}}`$. In terms of the polarization propagator in Eq. (11) one would get an infinite sum of diagrams such as those shown in Fig. 5.
We have already noted at the beginning of Section II A that, while for the two-body Green’s function $`G^{\text{ph}}`$ one can introduce an integral equation, this is not in general possible for the polarization propagator. It becomes possible when one approximates the irreducible vertex function $`\mathrm{\Gamma }^{13}`$ with the direct matrix elements of the interaction. In that case, in an infinite system one gets a simple algebraic equation whose solution, for the polarization propagators in Eq. (2) and the interaction in Eq. (2), is readily found to be
$$\mathrm{\Pi }_\text{X}^{\text{ring}}(q,\omega )=\frac{\mathrm{\Pi }^{(0)}(q,\omega )}{1\mathrm{\Pi }_\text{X}^{(1)\text{d}}(q,\omega )/\mathrm{\Pi }^{(0)}(q,\omega )},$$
(82)
where $`\mathrm{\Pi }_\text{X}^{(1)\text{d}}`$ represents the first-order direct polarization propagator:
$`\mathrm{\Pi }_{\text{L};I=0(1)}^{(1)\text{d}}(q,\omega )`$ $`=`$ $`\mathrm{\Pi }^{(0)}(q,\omega )4V_{0(\tau )}(q)\mathrm{\Pi }^{(0)}(q,\omega ),`$ (84)
$`\mathrm{\Pi }_{\text{T};I=0(1)}^{(1)\text{d}}(q,\omega )`$ $`=`$ $`\mathrm{\Pi }^{(0)}(q,\omega )4[V_{\sigma (\sigma \tau )}(q)V_{t(t\tau )}(q)]\mathrm{\Pi }^{(0)}(q,\omega ).`$ (85)
The effect of the exchange diagrams is often included through an effective zero-range interaction, calculated by taking the limit $`q0`$ of the first-order exchange contribution and rewriting it as an effective first-order direct term . Exact calculations, however, show that extrapolating this approximation to finite transferred momenta is not always reliable .
A more advanced approximation scheme is given by the continued fraction (CF) expansion . At infinite order the CF expansion exactly corresponds to the summation of the perturbative series, so that it is not any easier to calculate than the exact expression. However, when truncated at finite order, not only does it reproduce the standard perturbative series at the same order, but in addition it yields an estimate for each one of the infinite number of higher-order contributions. Regrettably no general methods are available to predict the convergence of the CF expansion, the only reliable test being to compare the results at successive orders.
On the other hand, one should note that for zero-range forces the first-order CF expansion already gives the exact (albeit trivial) result, making one hope that the short-range nature of the nuclear interactions allows for a fast convergence. Indeed, all available calculations have been performed truncating the CF expansion at first order . Here, as anticipated, we shall test the convergence up to second order.
The CF formalism for the polarization propagator is developed in Ref. for the case of Tamm-Dancoff correlations and extended in Ref. to the full RPA. Instead of following the rather involved formal derivation given there, here we shall briefly sketch a sort of heuristic derivation of the CF expansion.
Let us assume that we want to build a CF-like expansion for the polarization propagator, according to the pattern
$$\mathrm{\Pi }^{\text{RPA}}=\frac{\mathrm{\Pi }^{(0)}}{1A{\displaystyle \frac{B}{1C{\displaystyle \frac{D}{1\mathrm{}}}}}}.$$
(86)
We have said that the CF approach at $`n`$-th order exactly corresponds to the perturbative series at the same order and then it approximates the higher orders. Thus, if we want to approximate the exact RPA propagator at first order in CF (for sake of illustration we drop spin-isospin indices),
$$\mathrm{\Pi }^{\text{RPA}}=\underset{n=0}{\overset{\mathrm{}}{}}\mathrm{\Pi }^{(n)},$$
(87)
we can rather naturally set
$$\mathrm{\Pi }^{(n)}\mathrm{\Pi }^{(0)}\left[\frac{\mathrm{\Pi }^{(1)}}{\mathrm{\Pi }^{(0)}}\right]^n.$$
(88)
In Eq. (88) $`\mathrm{\Pi }^{(1)}\mathrm{\Pi }^{(0)}4V\mathrm{\Pi }^{(0)}+\mathrm{\Pi }^{(1)\text{ex}}`$ is the sum of the direct and exchange first-order terms of RPA — since this yields the correct expression for the direct terms. With the approximation in Eq. (88) the summation is trivial, yielding
$$\mathrm{\Pi }_{\text{CF1}}^{\text{RPA}}=\frac{\mathrm{\Pi }^{(0)}}{1\mathrm{\Pi }^{(1)}/\mathrm{\Pi }^{(0)}}=\frac{\mathrm{\Pi }^{(0)}}{14V\mathrm{\Pi }^{(0)}\mathrm{\Pi }^{(1)\text{ex}}/\mathrm{\Pi }^{(0)}}.$$
(89)
We could then add in the denominator of the above expression the exact second-order term, $`\mathrm{\Pi }^{(2)}`$, after subtracting its approximate estimate given by the first-order CF expansion, $`[\mathrm{\Pi }^{(1)}]^2/\mathrm{\Pi }^{(0)}`$. We would thus obtain
$`\mathrm{\Pi }_{\text{CF2}}^{\text{RPA}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Pi }^{(0)}}{1\mathrm{\Pi }^{(1)}/\mathrm{\Pi }^{(0)}\{\mathrm{\Pi }^{(2)}/\mathrm{\Pi }^{(0)}[\mathrm{\Pi }^{(1)}/\mathrm{\Pi }^{(0)}]^2\}}}`$ (90)
$`=`$ $`{\displaystyle \frac{\mathrm{\Pi }^{(0)}}{14V\mathrm{\Pi }^{(0)}\mathrm{\Pi }^{(1)\text{ex}}/\mathrm{\Pi }^{(0)}\{\mathrm{\Pi }^{(2)\text{ex}}/\mathrm{\Pi }^{(0)}[\mathrm{\Pi }^{(1)\text{ex}}/\mathrm{\Pi }^{(0)}]^2\}}}.`$ (91)
It is easily deduced from Eq. (91) that the third-order term is approximated as $`\mathrm{\Pi }^{(3)}\mathrm{\Pi }^{(1)}\{2\mathrm{\Pi }^{(2)}/\mathrm{\Pi }^{(0)}[\mathrm{\Pi }^{(1)}/\mathrm{\Pi }^{(0)}]^2\}`$. Then, going ahead in a CF-style expansion we would guess for the exact RPA propagator the following expression:
$$\mathrm{\Pi }^{\text{RPA}}=\frac{\mathrm{\Pi }^{(0)}}{1\mathrm{\Pi }^{(1)}/\mathrm{\Pi }^{(0)}{\displaystyle \frac{\mathrm{\Pi }^{(2)\text{ex}}/\mathrm{\Pi }^{(0)}[\mathrm{\Pi }^{(1)\text{ex}}/\mathrm{\Pi }^{(0)}]^2}{1{\displaystyle \frac{\mathrm{\Pi }^{(3)\text{ex}}/\mathrm{\Pi }^{(0)}+[\mathrm{\Pi }^{(1)\text{ex}}/\mathrm{\Pi }^{(0)}]^32[\mathrm{\Pi }^{(1)\text{ex}}/\mathrm{\Pi }^{(0)}][\mathrm{\Pi }^{(2)\text{ex}}/\mathrm{\Pi }^{(0)}]}{\mathrm{\Pi }^{(2)\text{ex}}/\mathrm{\Pi }^{(0)}[\mathrm{\Pi }^{(1)\text{ex}}/\mathrm{\Pi }^{(0)}]^2}}\mathrm{}}}}.$$
(92)
This is the expression that one would get from the formalism of Refs. when the expansion up to third order is worked out. Note that we did not assume any specific scheme (either Tamm-Dancoff or RPA) in this heuristic derivation.
Thus, following Eq. (82), we can write
$$\mathrm{\Pi }_\text{X}^{\text{RPA}}=\frac{\mathrm{\Pi }^{(0)}}{1\mathrm{\Pi }_\text{X}^{(1)\text{d}}/\mathrm{\Pi }^{(0)}\mathrm{\Pi }_\text{X}^{(1)\text{ex}}/\mathrm{\Pi }^{(0)}{\displaystyle \frac{\mathrm{\Pi }_\text{X}^{(2)\text{ex}}/\mathrm{\Pi }^{(0)}\left[\mathrm{\Pi }_\text{X}^{(1)\text{ex}}/\mathrm{\Pi }^{(0)}\right]^2}{1\mathrm{}}}},$$
(93)
where $`\mathrm{\Pi }_\text{X}^{(1)\text{d}}`$ has been defined in Eq. (82). Clearly, a truncation at $`n`$-th order would require the calculation of the exchange contributions up to that order. Exploiting Eq. (2) these can be cast in the form
$`\mathrm{\Pi }_{\text{L};I}^{(n)\text{ex}}(q,\omega )`$ $`=`$ $`\text{tr}[\widehat{O}_\text{L}^I\widehat{\mathrm{\Pi }}^{(n)\text{ex}}(q,\omega )\widehat{O}_\text{L}^I]={\displaystyle \underset{\alpha _i}{}}C_{\text{L};I}^{\alpha _1\mathrm{}\alpha _n}\mathrm{\Pi }_{\alpha _1\mathrm{}\alpha _n}^{(n)\text{ex}}(q,\omega ),`$ (95)
$`\mathrm{\Pi }_{\text{T};I}^{(n)\text{ex}}(q,\omega )`$ $`=`$ $`{\displaystyle \underset{ij}{}}\mathrm{\Lambda }_{ji}\text{tr}[\widehat{O}_{\text{T};i}^I\widehat{\mathrm{\Pi }}^{(n)\text{ex}}(q,\omega )\widehat{O}_{\text{T};j}^I]={\displaystyle \underset{\alpha _i}{}}C_{\text{T};I}^{\alpha _1\mathrm{}\alpha _n}\mathrm{\Pi }_{\alpha _1\mathrm{}\alpha _n}^{(n)\text{ex}}(q,\omega ),`$ (96)
where the indices $`\alpha _i`$ run over all the spin-isospin channels and the spin-isospin factors are absorbed into the coefficients $`C_\text{X}^{\alpha _1\mathrm{}\alpha _n}C_\text{X}^{(\alpha _1)}C_\text{X}^{(\alpha _2)}\mathrm{}C_\text{X}^{(\alpha _n)}`$ (see Table I).
Moreover the “elementary” exchange contribution $`\mathrm{\Pi }_{\alpha _1\mathrm{}\alpha _n}^{(n)\text{ex}}`$ containing $`n`$ interaction lines $`V_{\alpha _1}`$$`V_{\alpha _n}`$, namely<sup>§</sup><sup>§</sup>§ The following formulae are valid for non-tensor interactions; the treatment of the tensor terms is slightly more complex and it is given in Appendix C.
$`\mathrm{\Pi }_{\alpha _1\mathrm{}\alpha _n}^{(n)\text{ex}}(q,\omega )`$ $`=`$ $`i^{n+1}{\displaystyle \frac{d^4K_1}{(2\pi )^4}\mathrm{}\frac{d^4K_{n+1}}{(2\pi )^4}G^{(0)}(K_1)G^{(0)}(K_1+Q)V_{\alpha _1}(𝒌_1𝒌_2)\mathrm{}}`$ (98)
$`\mathrm{}V_{\alpha _n}(𝒌_n𝒌_{n+1})G^{(0)}(K_{n+1})G^{(0)}(K_{n+1}+Q)`$
$`=`$ $`(1)^n{\displaystyle \frac{d𝒌_1}{(2\pi )^3}\mathrm{}\frac{d𝒌_{n+1}}{(2\pi )^3}G_{\text{ph}}^{(0)}(𝒌_1,𝒒;\omega )V_{\alpha _1}(𝒌_1𝒌_2)\mathrm{}}`$ (100)
$`\mathrm{}V_{\alpha _n}(𝒌_n𝒌_{n+1})G_{\text{ph}}^{(0)}(𝒌_{n+1},𝒒;\omega )`$
have been introduced. With the definition of $`G_{\text{ph}}^{(0)}`$ given in Eq. (38) and by a suitable change of integration variables one can eliminate all of the $`\theta `$-functions that contain angular integration variables, leaving a multiple integral with the following general structure:
$`\mathrm{\Pi }_{\alpha _1\mathrm{}\alpha _n}^{(n)\text{ex}}(q,\omega )=(1)^n{\displaystyle \frac{d𝒌_1}{(2\pi )^3}\theta (k_Fk_1)\mathrm{}\frac{d𝒌_{n+1}}{(2\pi )^3}\theta (k_Fk_{n+1})}`$ (101)
$`\times [{\displaystyle \frac{1}{\omega ϵ_{𝒌_1+𝒒}+ϵ_{𝒌_1}+i\eta _\omega }}V_{\alpha _1}(𝒌_1𝒌_2)\mathrm{}V_{\alpha _n}(𝒌_n𝒌_{n+1}){\displaystyle \frac{1}{\omega ϵ_{𝒌_{n+1}+𝒒}+ϵ_{𝒌_{n+1}}+i\eta _\omega }}`$ (102)
$`+{\displaystyle }(\omega \omega )].`$ (103)
In Eq. (103), $`(\omega \omega )`$ stands for the sum of all the terms generated according to the following rules:
* Take all of the terms obtained by substituting $`\omega \omega `$ in one energy denominator in the second line of Eq. (103); then add the contribution obtained by performing the same substitution in two energy denominators and so on up to when the replacement $`\omega \omega `$ has been performed in all the $`n+1`$ denominators;
* Every time $`(\omega ϵ_{𝒌_i+𝒒}+ϵ_{𝒌_i}+i\eta _\omega )^1`$ is replaced with $`(\omega ϵ_{𝒌_i+𝒒}+ϵ_{𝒌_i}i\eta _\omega )^1`$ then replace $`𝒌_i`$ with $`𝒌_i𝒒`$ in the potential.
The number of integrations can be reduced by noticing that the azimuthal angles are contained only in the potential functions $`V_{\alpha _i}`$. For typical potentials this integration can be done analytically, hence it is convenient to introduce a new function representing the azimuthal integral of the potential. To this end, define the new variables:
$`|𝒌𝒌^{}|`$ $`=`$ $`\sqrt{k^2+k_{}^{}{}_{}{}^{2}2kk^{}[\mathrm{cos}\theta \mathrm{cos}\theta ^{}+\mathrm{sin}\theta \mathrm{sin}\theta ^{}\mathrm{cos}(\phi \phi ^{})]}`$ (104)
$`=`$ $`\sqrt{k^2+k_{}^{}{}_{}{}^{2}2[yy^{}+\sqrt{k^2y^2}\sqrt{k_{}^{}{}_{}{}^{2}y_{}^{}{}_{}{}^{2}}\mathrm{cos}(\phi \phi ^{})]}`$ (105)
$`=`$ $`\sqrt{x+x^{}2\sqrt{x}\sqrt{x^{}}\mathrm{cos}(\phi \phi ^{})+(yy^{})^2},`$ (106)
where $`yk\mathrm{cos}\theta `$ and $`xk^2y^2`$. Then, one can introduce
$$W_\alpha (x,y;x^{},y^{})=_0^{2\pi }\frac{d\phi }{2\pi }V_\alpha (𝒌𝒌^{})=W_\alpha (x^{},y^{};x,y)$$
(107)
and rewrite Eq. (103) as
$`\mathrm{\Pi }_{\alpha _1\mathrm{}\alpha _n}^{(n)\text{ex}}(q,\omega )`$ $`=`$ $`(1)^n\left({\displaystyle \frac{m_N}{q}}\right)^{n+1}\left({\displaystyle \frac{k_F}{2\pi }}\right)^{2n+2}{\displaystyle _1^1}𝑑y_1{\displaystyle \frac{1}{2}}{\displaystyle _0^{1y_1^2}}𝑑x_1\mathrm{}{\displaystyle _1^1}𝑑y_{n+1}{\displaystyle \frac{1}{2}}{\displaystyle _0^{1y_{n+1}^2}}𝑑x_{n+1}`$ (110)
$`\times {\displaystyle \frac{1}{\psi y_1+i\eta _\omega }}W_{\alpha _1}(x_1,y_1;x_2,y_2)\mathrm{}W_{\alpha _n}(x_n,y_n;x_{n+1},y_{n+1}){\displaystyle \frac{1}{\psi y_{n+1}+i\eta _\omega }}`$
$`+{\displaystyle (\omega \omega )}.`$
For $`n=1`$ one has
$`\mathrm{\Pi }_\alpha ^{(1)\text{ex}}(q,\omega )`$ $`=`$ $`\left({\displaystyle \frac{m_N}{q}}\right)^2{\displaystyle \frac{k_F^4}{(2\pi )^4}}{\displaystyle _1^1}𝑑y{\displaystyle \frac{1}{2}}{\displaystyle _0^{1y^2}}𝑑x{\displaystyle _1^1}𝑑y^{}{\displaystyle \frac{1}{2}}{\displaystyle _0^{1y_{}^{}{}_{}{}^{2}}}𝑑x^{}`$ (113)
$`\times {\displaystyle \frac{1}{\psi y+i\eta _\omega }}W_\alpha (x,y;x^{},y^{}){\displaystyle \frac{1}{\psi y^{}+i\eta _\omega }}`$
$`+{\displaystyle (\omega \omega )}`$
$`=`$ $`\left({\displaystyle \frac{m_N}{q}}\right)^2{\displaystyle \frac{k_F^4}{(2\pi )^4}}\left[𝒬_\alpha ^{(1)}(0,\psi )𝒬_\alpha ^{(1)}(\overline{q},\psi )+𝒬_\alpha ^{(1)}(0,\psi +\overline{q})𝒬_\alpha ^{(1)}(\overline{q},\psi +\overline{q})\right],`$ (114)
where
$$𝒬_\alpha ^{(1)}(\overline{q},\psi )=2_1^1𝑑y\frac{1}{\psi y+i\eta _\omega }_1^1𝑑y^{}W_{\alpha }^{}{}_{}{}^{\prime \prime }(y,y^{};\overline{q})\frac{1}{yy^{}+\overline{q}}$$
(116)
and
$$W_\alpha ^{\prime \prime }(y,y^{};\overline{q})=\frac{1}{2}_0^{1y^2}𝑑x\frac{1}{2}_0^{1y_{}^{}{}_{}{}^{2}}𝑑x^{}W_\alpha (x,y+\overline{q};x^{},y^{}).$$
(117)
Note that in getting to Eq. (116) use has been made of the Poincaré–Bertrand theorem . For the potential in Eq. (*II A) $`W_\alpha ^{\prime \prime }`$ can be calculated analytically (see Appendix D), so that the calculation of the first-order exchange contribution to the polarization propagator is reduced to the numerical evaluation of two-dimensional integrals for the real part and of one-dimensional integrals for the imaginary part.
For $`n=2`$ one has
$`\mathrm{\Pi }_{\alpha \alpha ^{}}^{(2)\text{ex}}(q,\omega )=\left({\displaystyle \frac{m_N}{q}}\right)^3{\displaystyle \frac{k_F^6}{(2\pi )^6}}{\displaystyle _1^1}𝑑y_1{\displaystyle \frac{1}{2}}{\displaystyle _0^{1y_1^2}}𝑑x_1{\displaystyle _1^1}𝑑y_2{\displaystyle \frac{1}{2}}{\displaystyle _0^{1y_2^2}}𝑑x_2{\displaystyle _1^1}𝑑y_3{\displaystyle \frac{1}{2}}{\displaystyle _0^{1y_3^2}}𝑑x_3`$ (118)
$`\times {\displaystyle \frac{1}{\psi y_1+i\eta _\omega }}W_\alpha (x_1,y_1;x_2,y_2){\displaystyle \frac{1}{\psi y_2+i\eta _\omega }}W_\alpha (x_2,y_2;x_3,y_3){\displaystyle \frac{1}{\psi y_3+i\eta _\omega }}`$ (119)
$`+{\displaystyle (\omega \omega )}`$ (120)
$`=\left({\displaystyle \frac{m_N}{q}}\right)^3{\displaystyle \frac{k_F^6}{(2\pi )^6}}[𝒬_{\alpha \alpha ^{}}^{(2)}(0,0;\psi )𝒬_{\alpha \alpha ^{}}^{(2)}(0,\overline{q};\psi )𝒬_{\alpha \alpha ^{}}^{(2)}(\overline{q},0;\psi )+𝒬_{\alpha \alpha ^{}}^{(2)}(\overline{q},\overline{q};\psi )`$ (121)
$`𝒬_{\alpha \alpha ^{}}^{(2)}(0,0;\psi +\overline{q})+𝒬_{\alpha \alpha ^{}}^{(2)}(0,\overline{q};\psi +\overline{q})+𝒬_{\alpha \alpha ^{}}^{(2)}(\overline{q},0;\psi +\overline{q})𝒬_{\alpha \alpha ^{}}^{(2)}(\overline{q},\overline{q};\psi +\overline{q})],`$ (122)
(123)
where
$$𝒬_{\alpha \alpha ^{}}^{(2)}(\overline{q}_1,\overline{q}_2;\psi )=_1^1𝑑y\frac{1}{2}_0^{1y^2}𝑑x𝒢_\alpha (x,y+\overline{q}_1;\psi +\overline{q}_1)\frac{1}{\psi y+i\eta _\omega }𝒢_\alpha ^{}(x,y+\overline{q}_2;\psi +\overline{q}_2)$$
(124)
and
$`𝒢_\alpha (x,y;\psi )`$ $`=`$ $`{\displaystyle _1^1}𝑑y^{}{\displaystyle \frac{1}{\psi y^{}+i\eta _\omega }}W_\alpha ^{}(x,y;y^{}),`$ (126)
$`W_\alpha ^{}(x,y;y^{})`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _0^{1y_{}^{}{}_{}{}^{2}}}𝑑x^{}W_\alpha (x,y;x^{},y^{}).`$ (127)
For the potential in Eq. (*II A) $`W_\alpha ^{}`$ can be calculated analytically (see Appendix D) and one is left with the numerical integration of Eqs. (124) and (126), so that the calculation of the second-order exchange contribution to the polarization propagator is effectively reduced to the numerical evaluation of at most three-dimensional integrals. Higher orders add a numerical two-dimensional integration for each additional interaction line, since, for a potential of the form in Eq. (*II A), only the azimuthal integration can be performed analytically for the interaction lines that do not close on the external vertices.
Finally we recall that the nucleon propagators can be dressed by the HF field, as explained in Section II D, by replacing $`\psi \psi ^{}`$ and $`m_Nm_N^{}`$, where $`\psi ^{}`$ and $`m_N^{}`$ have been defined in Eqs. (72) and (77), multiplying by the appropriate power of the normalization factor $`1+\omega /m_N+\mathrm{\Delta }^{(1)}/m_N`$ when relativistic kinematics are employed (see Eqs. (3) and (81)).
### F Effective particle-hole interaction
In order to assess the contributions to the nuclear responses arising from the various approximation schemes introduced so far and to show typical results, first of all we have to choose an effective interaction. This choice can be rather delicate as it may introduce uncontrolled uncertainties in the calculation. Here, however, we are not interested so much in comparisons with data, but rather with setting up working many-body schemes. For this purpose, we shall use the $`G`$-matrix based on the Bonn potential of Ref. , adapted to the quasielastic regime as in Ref. . Although the attraction provided in the scalar-isoscalar channel by this interaction is definitely too strong , it will serve our illustrative needs.
Two approaches to determine the effective ph interaction in the nuclear medium appear to be possible: one can either directly fix an effective potential by fitting some phenomenological properties or start with a bare nucleon-nucleon interaction and calculate the related $`G`$-matrix. Parameterizations of the ph interaction based upon the first procedure are generally only available at very low momentum transfers (in terms of Migdal-Landau parameters), and since we are probing relatively high momenta, we have resorted to using a $`G`$-matrix. We have chosen the one of Ref. , that, in our view, has the following appealing features: it is based upon a realistic boson-exchange potential; it accounts for the density dependence; and it includes (nonlocal) exchange contributions in the effective interaction, which are conveniently parameterized in terms of Yukawa functions.
A feature related to the effective inclusion of antisymmetrization effects is particularly interesting in order to test a specific widely employed approximation scheme, the so-called ring approximation, in which the exchange diagrams of the RPA series are dropped and their effect mimicked by adding to the direct interaction matrix elements an effective exchange contribution (see, e. g., Ref. ). Indeed, below we shall compare calculations employing the fully antisymmetrized formalism developed in the previous subsections using the direct part of the $`G`$-matrix, to those employing the ring approximation using the antisymmetrized effective interaction.
To facilitate the comparison with the original parameterization of Ref. , the potential is given here using the standard representation of Eq. (2) in spin and isospin (no spin-orbit contribution will be considered in the following), but employing different symbols for the momentum space potentials (and adding the tensor contributions in the exchange channel):
$`V(𝒌_f,𝒌_i;k_F)`$ $`=`$ $`F+F^{}𝝉_1𝝉_2+G𝝈_1𝝈_2+G^{}𝝈_1𝝈_2𝝉_1𝝉_2`$ (129)
$`+TS_{12}(\widehat{𝒒})+T^{}S_{12}(\widehat{𝒒})𝝉_1𝝉_2+HS_{12}(\widehat{𝑸})+H^{}S_{12}(\widehat{𝑸})𝝉_1𝝉_2,`$
where $`𝒒=𝒌_i𝒌_f`$, $`𝑸=𝒌_i+𝒌_f`$ ($`𝒌_i`$, $`𝒌_f`$ being the relative momenta in the initial and final states, respectively) and the coefficients are density and momentum dependent.
Before utilizing the interaction of Ref. in a calculation of quasielastic responses, a few issues have to be addressed .
a) The density dependence of the $`G`$-matrix is given in terms of density-dependent coupling constants, which is not very useful for applications to finite nuclei. Furthermore, the parameterization is fitted for 0.95 fm$`{}_{}{}^{1}<k_F<`$ 1.36 fm<sup>-1</sup>, which spans a range of densities down to roughly 1/3 of the central density. Extrapolation of this parameterization to lower densities (which is crucial for application to hadron scattering) gives unreasonable results. Thus, we have chosen to employ a linear $`\rho `$-dependence ($`V=V^{\text{ex}}+V^\rho \rho `$), which is considered a reasonable choice (see, e. g., Ref. ). In Fig. 7 one can see a comparison of the two parameterizations for the $`k_F`$-dependence of the effective interaction. It should be noted that most of the contribution to the quasielastic responses comes from densities where the two descriptions differ by only a few percent.
b) In order to obtain a local interaction at a fixed density, one can use the relation between $`q`$, $`Q`$ and $`k_i`$, i. e., $`Q=\sqrt{4k_i^2q^2}`$ and take for $`k_i`$ a suitably chosen average value, $`k_i`$; then, the only independent momentum is $`q`$. The authors of Ref. were interested in a potential for low excitation energy nuclear structure calculations and hence they assumed that the two nucleons in the initial state lie on the Fermi surface and so averaged over the relative angle, getting $`k_i0.7k_F`$. Clearly, in this case one has the constraint $`0<q1.4k_F`$. On the other hand, we are interested in the ph interaction in the quasielastic region where one nucleon in the initial state is below the Fermi sea, while the other can be well above it. A look at Fig. 7 shows that $`𝒌_i`$ is defined in terms of the particle and hole momenta as $`𝒌_i=(𝒑𝒉^{})/2=(𝒉𝒉^{}+𝒒)/2`$. Thus, at fixed $`𝒒`$ one should average $`k_i`$ over $`𝒉`$ and $`𝒉^{}`$, getting $`k_i\sqrt{6k_F^2/5+q^2}/2`$. Now $`k_i`$ grows with $`q`$, so that there are no longer constraints on $`q`$ and the exchange momentum turns out to be constant, $`Q=\sqrt{6/5}k_F`$. In Fig. 9 one can see the resulting interaction in the non-tensor channels.
c) The tensor channels are simpler, since in the parameterization of Ref. there is no explicit density dependence (Fig. 9). The coefficients of the exchange tensor operator, $`H`$ and $`H^{}`$, display a very mild density dependence induced by $`Q`$, which is completely negligible. The only drawback concerns the treatment of $`S_{12}(𝑸)`$: assuming that $`𝒒`$ and $`𝑸`$ are orthogonal, with some algebra one can show that $`S_{12}(\widehat{𝑸})=S_{12}(\widehat{𝒒})/2`$.
### G Testing the model
First of all we have to choose the Fermi momentum. Of course, one could easily perform a local density calculation to achieve a better description of finite nuclei. Here, for sake of illustration, we prefer to use the pure Fermi gas. The choice of $`k_F`$ can be made in several ways — here we shall choose an average value according to the formula (see, e. g., )
$$\overline{k}_F=\frac{1}{A}𝑑𝒓k_F(r)\rho (r),$$
(130)
where $`\rho (r)`$ is the empirical Fermi density distribution normalized to the number of nucleons and $`k_F(r)=[(3\pi /2)\rho (r)]^{1/3}`$. For <sup>12</sup>C one gets $`\overline{k}_F195`$ MeV/c and this is the value used in the calculations that follow.
Let us start with the HF response. In Fig. 10 we display the HF response of <sup>12</sup>C at $`q=300`$, 500 and 1000 MeV/c. As anticipated, Eq. (79) turns out to be a good approximation to the exact expression (54) (except on the borders of the response region, where the Fermi gas is anyway unrealistic). The HF correlations widen the response region and quench and harden the position of the quasielastic peak, as is well known. Note however that the short-range correlations, which are embodied in the effective interaction based on a $`G`$-matrix, reduce the amount of hardening that is observed in calculations based on the bare Bonn potential . Note also that the same level of accuracy is obtained using either non-relativistic or relativistic kinematics.
Before discussing the RPA results, we would like to test the convergence of the CF expansion. For this purpose, in Fig. 12 we compare the longitudinal RPA responses at first and second order in the CF expansion using a model one-boson-exchange interaction, $`V_\sigma (k)=𝝈_1𝝈_2g[m^2/(m^2+k^2)]`$ (the spin operators having the purpose of killing the direct (ring) contribution). For values of the coupling constant $`g`$ and of the boson mass $`m`$ typical of realistic nucleon-nucleon potentials one finds that the first- and second-order results match at the level of a few percent (in the left and middle panels of Fig. 12, the solid and dashed curves are actually indistinguishable). One has to go to very low boson masses (a few MeV) and, consequently to very high values of $`g`$ in order to find some discrepancies. To understand these results better, in Fig. 12 we display the modulus of the polarization propagator at first order, $`\mathrm{\Pi }^{(1)}`$ (dotted), at second order, $`\mathrm{\Pi }^{(2)}`$ (dashed) and the approximation to $`\mathrm{\Pi }^{(2)}`$ generated by the first-order CF expansion (see Section II E), $`\mathrm{\Pi }^{(2)\text{appr}}\mathrm{\Pi }_{}^{(1)}{}_{}{}^{2}/\mathrm{\Pi }^{(0)}`$ (solid). From inspection of the curves, it is clear that the first important element to guarantee good convergence is the range of the interaction. Indeed, for $`m=800`$ MeV (short-range), $`\mathrm{\Pi }^{(2)}`$ and $`\mathrm{\Pi }^{(2)\text{appr}}`$ practically coincide independent of the strength of the interaction. This, of course, should be expected, since for zero-range interactions the first-order CF expansion gives the exact result. For masses of the order of the pion mass one starts finding discrepancies between $`\mathrm{\Pi }^{(2)}`$ and $`\mathrm{\Pi }^{(2)\text{appr}}`$. However, for realistic values of the interaction strength the second-order contribution turns out to be one order of magnitude smaller than the first-order one, and thus these discrepancies have little effect on the full response functions (Fig. 12).
To understand these results it may be useful to compare the strength of the interactions employed here to that of one-pion-exchange, $`g_\pi m_\pi ^2f_\pi ^2/30.33`$ (in natural units). With the same units, the cases with $`m=100`$ MeV correspond to $`gm^2=0.26`$ and 0.65; those with $`m=800`$ MeV to $`gm^2=16.7`$ and 41.7; for $`m=1`$ and 10 MeV one has $`gm^2=1.3`$ and 0.65, respectively.
To summarize, from the left and middle panels of Fig. 12 one can understand that the validity of the CF expansion originates from the interplay between range and strength of the interaction. For short-range potentials where the conventional perturbative expansion may not converge, the CF technique yields a good approximation for the propagators at all orders; for long-range (on the nuclear scale) forces, the CF approximation is less accurate, but the relative weakness of the interaction already guarantees the convergence of the conventional perturbative expansion. One has to go to unreasonably low masses to find a situation where the interaction range is very long and $`\mathrm{\Pi }^{(1)}`$ and $`\mathrm{\Pi }^{(2)}`$ are of the same order (right panels in Fig. 12).
We can thus conclude that the calculations of nuclear response functions in the antisymmetrized RPA performed at first order in the CF expansion are indeed quite accurate. The same conclusion is also supported by calculations with a realistic effective interaction — such as the $`G`$-matrix parameterization discussed above — and including HF and relativistic kinematical effects. In fact, in Fig. 14 we show the RPA and BHF-RPA longitudinal responses of <sup>12</sup>C at $`q=300`$, 500 and 1000 MeV/c, using the full $`G`$-matrix introduced at the beginning of this section. Also for the full interaction, the discrepancies between the first- and second-order CF responses are too small to be displayed. They are at the level of fractions of percent everywhere, except for the case of the isoscalar channel at 300 MeV/c, where they rise to a few percent due to the closeness of a singularity in the propagator induced by the strongly attractive interaction. Indeed, as already mentioned, the scalar-isoscalar channel is (too) attractive As indicated by the energy position of the breathing modes; in other words, nuclear matter with such an interaction becomes unstable. and softens the quasielastic peak; the scalar-isovector one is repulsive and gives rise to a hardening. The effect of the HF correlations is the same as in the discussion of Fig. 10. In Fig. 14 the transverse response is displayed for the same conditions.
Finally, it is interesting and important to test the validity of the ring approximation — where exchange diagrams are not included — since this approximation has been widely used in the literature because of its simplicity. In this scheme, the effect of antisymmetrization is simulated by adding to the direct interaction matrix elements an effective exchange contribution (see, e. g., Ref. ). For details see also Ref. , where a prescription to determine the effective exchange momentum designed for use in the quasifree region has been given.
In Fig. 15 we display the ring and RPA responses of <sup>12</sup>C at $`q=500`$ MeV/c, using the $`G`$-matrix parameterization. It is apparent that the only channel where the ring approximation works reasonably well is the spin-isovector one, which, incidentally, is the dominant one in ($`e`$,$`e^{}`$) magnetic scattering; it is less accurate in all other channels, especially in the scalar-isoscalar one. The same considerations also apply when the HF mean field is included in the ring and RPA responses. Note that these results confirm those of Ref. , where a comparison of ring and RPA calculations had been done using a numerically rather involved finite nucleus formalism. Also in that calculation the $`G`$-matrix of Ref. had been employed.
## III Parity-violating electron scattering and axial responses
### A The asymmetry, the currents and the RFG responses
A new window on the inclusive nuclear responses that allows us to unravel aspects of nuclear and nucleon structure that are otherwise inaccessible to unpolarized probes is offered by parity-violating electron scattering from nuclei. See Ref. for a general review of the subject. Experiments of this type exploit longitudinally polarized electrons to measure the helicity asymmetry $`𝒜`$, defined as the difference between the inclusive nuclear scattering of right- and left-handed electrons divided by their sum, namely
$$𝒜=\frac{d^2\sigma ^+d^2\sigma ^{}}{d^2\sigma ^++d^2\sigma ^{}}.$$
(131)
$`𝒜`$ arises from the interference between the electromagnetic current that is purely vector (V) and the weak neutral current that has both vector and axial-vector (A) components. Diagrams for the associated amplitudes in leading order of the bosons exchanged (the photon and the $`Z^0`$) are displayed in Fig. 16. In this approximation Eq. (131) can be cast in the form
$$𝒜=𝒜_0\frac{v_\text{L}R_\text{L}^{\text{A}V}(q,\omega )+v_\text{T}R_\text{T}^{\text{A}V}(q,\omega )+v_\text{T}^{}R_\text{T}^{}^{\text{V}A}(q,\omega )}{v_\text{L}R_\text{L}(q,\omega )+v_\text{T}R_\text{T}(q,\omega )}.$$
(132)
The numerator is parity-violating (PV) and the denominator parity-conserving (PC). Here
$$v_\text{L}=\left(\frac{Q^2}{𝒒^2}\right)^2$$
(133)
$$v_\text{T}=\frac{1}{2}\left|\frac{Q^2}{𝒒^2}\right|+\mathrm{tan}^2\frac{\theta }{2}$$
(134)
and
$$v_\text{T}^{}=\mathrm{tan}\frac{\theta }{2}\sqrt{\left|\frac{Q^2}{𝒒^2}\right|+\mathrm{tan}^2\frac{\theta }{2}}$$
(135)
are the usual leptonic kinematical factors, $`\theta `$ is the electron scattering angle and, as before, $`Q^2=\omega ^2𝒒^2`$ is the spacelike four-momentum transferred from the electron to the nucleus.
In Eq. (131) the nuclear (and nucleon’s) structure are embedded in the electromagnetic PC nuclear responses $`R_\text{L}`$ and $`R_\text{T}`$ discussed in the previous sections, while their PV analogs $`R_\text{L}^{\text{A}V}`$, $`R_\text{T}^{\text{A}V}`$ and $`R_\text{T}^{}^{\text{V}A}`$ are discussed in this section. Here the first (second) index in the superscript refers to the vector (V) or axial (A) nature of the leptonic (hadronic) WNC. For brevity we shall often simply refer to these by their hadronic character, i.e., the L and T PV responses are called “vector” and the $`T^{}`$ response “axial”. Finally the scale of the asymmetry is set by the factor
$$𝒜_0=\frac{\sqrt{2}G_Fm_N^2}{\pi \alpha }\frac{|Q^2|}{4m_N^2}6.5\times 10^4\tau ,$$
(136)
which is defined in terms of the EM $`(\alpha )`$ and Fermi $`(G_F)`$ coupling constants. If there were no additional dependence on $`q`$ and $`\omega `$, then the expression for $`𝒜_0`$ would imply that the asymmetry grows with $`\tau =|Q^2|/4m_N^2`$ and hence it is not surprising that the first parity violation in electron scattering was observed at high energies at SLAC . On the other hand, it is also clear that selective processes such as elastic scattering do contain additional dependences on $`(q,\omega )`$ via form factors that may make measurements at large $`\tau `$ extremely difficult. In fact, only a very few have been performed to date. Even more challenging, but not impossible, are experiments whose goal is to disentangle in Eq. (132) the separate contributions of the PV responses $`R_\text{L}^{\text{A}V}`$, $`R_\text{T}^{\text{A}V}`$ and $`R_\text{T}^{}^{\text{V}A}`$.
In the investigation of the PV nuclear responses — these can assume positive as well as negative values — of central importance is the isospin decomposition of the hadronic four-current
$$\left(J_\mu \right)_{I,I_z}=\beta ^{(0)}\left(J_\mu \right)_{0,0}+\beta ^{(1)}\left(J_\mu \right)_{1,0}$$
(137)
into isoscalar ($`I`$=0) and isovector ($`I`$=1) components ($`\mu `$ is the Lorentz index). In the Standard Model at tree level the coefficients in Eq. (137) in the EM sector read
$$\beta _{\text{V},EM}^{(0)}=\beta _{\text{V},EM}^{(1)}=\frac{1}{2},$$
(138)
i.e. one has the usual responses $`R_\text{L}`$ and $`R_\text{T}`$. In the WNC sector, instead, they read
$`\beta _{\text{V},WNC}^{(0)}`$ $`=`$ $`2\mathrm{sin}^2\theta _W0.461`$ (139)
$`\beta _{\text{V},WNC}^{(1)}`$ $`=`$ $`12\mathrm{sin}^2\theta _W0.539`$ (140)
for the vector coupling and
$`\beta _\text{A}^{(0)}`$ $`=`$ $`0`$ (141)
$`\beta _\text{A}^{(1)}`$ $`=`$ $`1`$ (142)
for the axial one. The above results obtain with the following value of the weak mixing angle
$$\mathrm{sin}^2\theta _W=0.23055\pm 0.00041.$$
(143)
The importance of isospin becomes clear if one assumes it to be an exact symmetry in nuclei (which is of course only approximately true). Then for PV elastic scattering on spin zero, isospin zero nuclei, Eq. (132) reduces to
$$\frac{𝒜}{𝒜_0}=2\mathrm{sin}^2\theta _W$$
(144)
(see Eq. (156) for the coefficient of the axial leptonic current), which suggests using PV experiments as a tool for testing the Standard Model in the low-energy regime (see, for example, Ref. ).
To this point we have been assuming that the strangeness content in the nucleon or nucleus is negligible. If this is not the case, then it has been realized that PV electron scattering can be invaluable in exploring this aspect of hadronic structure (again, see Ref. for a review that contains discussion of this issue). Indeed, when strangeness content is taken into account Eq. (144) is modified as follows
$$\frac{𝒜}{𝒜_0}=2\mathrm{sin}^2\theta _W+\frac{G_E^{(s)}}{G_E^{(0)}}.$$
(145)
In the above $`G_E^{(0)}=G_{E_p}+G_{E_n}`$ and $`G_E^{(s)}`$ are the electric isoscalar EM and strange form factors of the nucleon (the indices $`p`$ and $`n`$ refer to the proton and the neutron, respectively). Formula (145) will be exploited to extract $`G_E^{(s)}`$ in an experiment planned at CEBAF involving elastic PV scattering from <sup>4</sup>He.
A further clue to the strangeness content may be seen in studying PV elastic scattering from the proton: in this case Eq. (132) becomes
$$\frac{𝒜}{𝒜_0}=\frac{\epsilon G_{E_p}\stackrel{~}{G}_{E_p}+\tau G_{M_p}\stackrel{~}{G}_{M_p}+\delta G_{M_p}\stackrel{~}{G}_{A_p}}{\epsilon G_{E_p}^2+\tau G_{M_p}^2},$$
(146)
where
$$\epsilon =\frac{1}{1+2(1+\tau )\mathrm{tan}^2(\theta /2)}$$
(147)
and
$$\delta =(14\mathrm{sin}^2\theta _W)\sqrt{\tau (1+\tau )(1\epsilon ^2)}.$$
(148)
In Eq. (146) $`\stackrel{~}{G}_{E_p}`$, $`\stackrel{~}{G}_{M_p}`$ and $`\stackrel{~}{G}_{A_p}`$ are the electric, magnetic and axial weak form factors of the proton, that read
$`\stackrel{~}{G}_{E_p,M_p}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[\beta _{\text{V},WNC}^{(0)}+\beta _{\text{V},WNC}^{(1)}\right]G_{E_p,M_p}+\left[\beta _{\text{V},WNC}^{(0)}\beta _{\text{V},WNC}^{(1)}\right]G_{E_n,M_n}`$ (149)
$`=`$ $`{\displaystyle \frac{1}{2}}(14\mathrm{sin}^2\theta _W)G_{E_p,M_p}{\displaystyle \frac{1}{2}}G_{E_n,M_n}{\displaystyle \frac{1}{2}}G_{E,M}^{(s)}`$ (150)
and
$$\stackrel{~}{G}_{A_p}=\frac{1}{2}G_{A_p}\frac{1}{2}G_{A_n}\frac{1}{2}G_A^{(s)}.$$
(151)
Equation (146) has already been exploited in experiments performed at Bates (SAMPLE) and CEBAF (HAPPEX) (see Refs. and , respectively) for unraveling the strangeness content of the proton. Figure 17 may help in understanding the magnitudes of the quantities given above. There one infers that at backward scattering angles (SAMPLE) it is mainly the magnetic strangeness that is measured, however with some contamination arising from the axial contribution, whereas at forward angles (HAPPEX) it is a mixture of electric and magnetic strangeness that is observed, the axial contribution being totally negligible in that case.
Although the above discussions of strangeness content are focused on the responses of the nucleon rather than of the nucleus, in fact it is also relevant for the latter, since the nuclear responses measured in the quasielastic regime are indeed affected by the nucleon’s isoscalar form factors and these in turn are affected by strangeness. In fact this is an excellent example of how the theoretical predictions in nuclear many-body and particle physics are interrelated.
In particular, with regard to the responses of the nucleus, PV experiments offer the opportunity of
* disentangling the isoscalar and isovector contributions to $`R_\text{L}^{\text{A}V}`$, $`R_\text{T}^{\text{A}V}`$ and $`R_\text{T}^{}^{\text{V}A}`$,
* exploring the Coulomb sum rule separately in the isoscalar and isovector channels (see Section IV),
* measuring the neutron distribution in nuclei (see Ref. ; not discussed here),
* investigating the nuclear axial response especially in the $`\mathrm{\Delta }`$ region (also not discussed in the present work)
and
* unambiguously revealing the role of the pion in nuclear excitations through the (possible) existence of a zero in the frequency behavior of the asymmetry.
To see how this occurs let us first split the PC and PV responses into their isospin components (we are neglecting the strangeness content at this point) according to
$`R_{\text{L},\text{T}}=\beta _{\text{V},EM}^{(0)}R_{\text{L},\text{T}}^{I=0}+\beta _{\text{V},EM}^{(1)}R_{\text{L},\text{T}}^{I=1}`$ (152)
$`R_{\text{L},\text{T}}^{\text{A}V}=\beta _{\text{V},WNC}^{(0)}R_{\text{L},\text{T}}^{I=0}+\beta _{\text{V},WNC}^{(1)}R_{\text{L},\text{T}}^{I=1}`$ (153)
for the vector channel. The PV axial channel is purely isovector at tree level. Next let us focus on the RFG model where the longitudinal isoscalar response is essentially proportional to $`\left(G_E^{(0)}\right)^2=\left(G_{E_p}+G_{E_n}\right)^2`$ and the isovector one to $`\left(G_E^{(1)}\right)^2=\left(G_{E_p}G_{E_n}\right)^2`$. Since $`G_{E_n}`$ is small, especially at low $`\tau `$, it follows that
$$R_\text{L}^{I=0}R_\text{L}^{I=1}.$$
(154)
But then it is clear that the RFG PV longitudinal response is almost vanishing because of the opposite sign and approximately equal magnitude of the coefficients in Eqs. (139) and (140). This dramatic consequence of the Standard Model is displayed in Fig. 18, where the five responses entering in the definition of the asymmetry in Eq. (132) are shown for $`|𝒒|`$=300, 500 and 2000 MeV/c together with $`𝒜`$.
Beyond the fact that $`R_\text{L}^{\text{A}V}`$ is very small, one also observes in the figure that
* a similar cancellation does not occur in the transverse channel, where in fact the isoscalar and the isovector responses are quite different because now $`(G_M^{(0)})^2(\mu _p+\mu _n)^20.77`$ and $`(G_M^{(1)})^2(\mu _p\mu _n)^222.1`$;
* $`R_\text{T}^{\text{A}V}`$ and $`R_\text{T}^{}^{\text{V}A}`$ are negative, this being related to the sign of the axial coefficient of the leptonic WNC
$$j_\mu (k^{},s^{};k,s)^{WNC}=\overline{u}(k^{},s^{})(a_V\gamma _\mu +a_A\gamma _5\gamma _\mu )u(k,s),$$
(155)
where $`u(k,s)`$ is the Dirac spinor of the electron. Indeed, according to the Standard Model,
$$a_V=(14\mathrm{sin}^2\theta _W)0.092\text{and}a_A=1;$$
(156)
* as a consequence the asymmetry is negative as well, reflecting the left-handed nature of the weak interaction;
* the asymmetry, as previously mentioned, grows with $`\tau `$ and $`\theta `$.
How do interactions among the constituents of the RFG modify the above predictions?
### B The role of the pion and other mesons
In answering the last question of the previous section we extend our model from a strict RFG framework to one where pions are also included , because then we can at least approximately preserve the two major requirements of Lorentz covariance and gauge invariance. Indeed the RFG cross section is built (apart from overall kinematical factors) from the contraction of the leptonic and hadronic Lorentz tensors and is therefore a relativistic invariant, although the partition into longitudinal and transverse responses depends, of course, upon the reference system. Moreover the single-nucleon four-current entering into the RFG nuclear tensor is conserved and hence the non-interacting RFG is gauge invariant.
When the nucleon-nucleon interaction carried by the pion is switched on it is not obvious that the two above mentioned properties are retained. Indeed, when the correlations and meson exchange currents (MEC) associated with the one pion exchange potential (OPEP) are introduced, as shown in Ref. the Lorentz covariance and gauge invariance are violated (however, only slightly so) due to the following approximations that are usually made:
i) the pion propagator is assumed to be static,
ii) and a non-relativistic expansion of the two-body currents is performed .
In order to achieve a treatment of forces and currents that is as consistent as possible we first limit ourselves to the study of diagrams with only one pionic line, namely, we work in the first perturbative order in the N-N interaction. The correlation diagrams to be evaluated in this scheme are the so-called self-energy and exchange contributions, that when iterated to infinite order generate the previously discussed HF and RPA series, respectively. Note that the tadpole and ring diagrams vanish due to the spin and isospin structure of the OPEP. Concerning MEC, three contributions occur, the pion-in-flight term, the contact term and the one associated with the $`\mathrm{\Delta }`$. A direct comparison with the exact relativistic calculation shows that the non-relativistic expansion of Ref. is indeed quite accurate up to momentum transfers of the order of 1 GeV/c.
The outcome of this is that sizable pionic contributions to the EM longitudinal (spin scalar, $`\sigma =0`$) and transverse (spin vector, $`\sigma =1`$) nuclear responses are found. In both cases the correlation effects produce a hardening of the responses, that is, a shift of the strength to higher $`\omega `$. The PV longitudinal and transverse correlated responses are simply obtained from the EM ones through the isospin rotations implied by the structure of the WNC discussed in the previous subsection; the axial-vector response will be treated separately in Section III C.
The main points emerging from this analysis are:
* In isospace the contribution of the self-energy diagram to the charge response is almost equally split between isoscalar ($`I=0`$) and isovector ($`I=1`$) components. The latter, on the other hand, is of course overwhelming in the transverse response, due to the dominance of the isovector magnetic moment. In contrast, in the case of the pionic force the $`I=0`$ part of the exchange diagram turns out to be three times as large as the $`I=1`$ one in the charge response and this imbalance, that becomes even stronger in higher orders of perturbation theory, is further strengthened by the difference between the isoscalar and isovector form factors. The isoscalar dominance of the pionic exchange correlations has dramatic consequences for the PV longitudinal response function, as may be seen in Fig. 19, where this response is displayed as a function of $`\omega `$ with and without pionic correlations.
A physical interpretation of why $`R_\text{L}^{\text{A}V}`$ is small in the independent-particle model and why isospin-correlations are so important in determining its ultimate size follows from the expression for the observables in a language that explicitly refers to neutrons and protons rather than employing isospin labeling. Indeed, by inspecting Fig. 20, where the diagrams describing both the EM and PV longitudinal responses for a free system are displayed, one easily understands why in the non-interacting case the EM longitudinal response turns out to be substantial: both of its vertices can in fact be large, as they involve the coupling of a longitudinal photon to a proton. In contrast, one of the vertices entering in the non-interacting PV response is always small, since either the longitudinal coupling of a photon to a neutron or of a $`Z^0`$ to a proton is involved. This last fact is often phrased by saying that the $`\gamma `$ is blind to neutrons and the $`Z^0`$ is blind to protons (i.e., in the longitudinal channel). To quantify the meaning of “large” and “small” we note that typically the $`\gamma n`$ coupling is about 1/10 that of the $`\gamma p`$, and likewise the $`Z^0p`$ coupling is about 1/10 that of the $`Z^0n`$.
The exchange correlations corresponding to the exchange of an isovector charged meson between the particle and hole convert a neutron (proton) into a proton (neutron), and thus give rise to a diagram where both couplings are large — hence the crucial role of such isovector correlations in determining $`R_\text{L}^{\text{A}V}`$.
The above arguments clearly do not apply to the transverse response both because in this case the $`I=0`$ channel is much weaker than the $`I=1`$ channel, being essentially proportional to the squares of the very different isoscalar and isovector magnetic moments of the nucleon, and furthermore because protons and neutrons can both couple strongly to photons via their (comparable) magnetic moments. Being also essentially isovector, the axial-vector response likewise does not display the sensitivity expected for the longitudinal response.
* While the tensor component of the OPEP never contributes to the self-energy in a translationally invariant system, the exchange diagram gets a tensor contribution, but only in the transverse channel and mostly via the backward-going graphs. This implies a different role for the pionic force in the two EM responses, a finding that should be tested against experiment.
* When we extend our analysis from first to infinite order of perturbation theory, thus generating the HF and RPA responses, the results do not change substantially. Although the pionic interaction is strong, it nevertheless therefore appears that for quasielastic kinematics its effects are reasonably small at not too small $`q`$, thus rendering perturbation theory quite accurate already at the lowest order, at least for the classes of diagrams studied here.
* The contribution of the central part of the pionic interaction to the exchange diagram stems largely from the $`\delta `$-force and not from the finite-range one. Notably for pointlike nucleons the continuity equation is obeyed by the OPEP and by the related pionic MEC, but is not affected by its $`\delta `$ component. However in keeping with the usual approach taken in studies of pionic effects, we include a $`\pi `$NN vertex function $`\mathrm{\Gamma }_\pi `$, whose scale is set by a mass parameter $`\mathrm{\Lambda }_\pi `$, to smear out the $`\delta `$-piece of OPEP. As a consequence the continuity equation is modified by the presence of $`\mathrm{\Gamma }_\pi `$ and to restore its validity additional MEC should be introduced. Significantly these counterterms affect only the pion-in-flight current, which is tiny, and therefore are quantitatively negligible.
* In contrast to the exchange diagram, the self-energy term gets a contribution only from momentum-dependent forces, and therefore an unmodified pionic $`\delta `$-interaction does not contribute in this channel. In fact, the contribution coming from the $`\delta `$-function in OPEP, when modified by the $`\pi `$NN vertex form factor, contains an effective momentum-dependence and so is nonzero.
* Most remarkably when the interactions carried by heavier mesons (namely $`\rho `$, $`\sigma `$ and $`\omega `$) are switched on, the previous conclusions on the PV responses (in particular the dramatic enhancement of $`R_\text{L}^{\text{A}V}`$) are not substantially changed, as illustrated in Ref. , thus showing the central role played by the pion in the nuclear dynamics for quasielastic kinematics. This is accomplished, however, through rather subtle aspects of the nuclear many-body problem, where the interference between the pion and the other mesons turns out to be crucial.
### C The axial response and the asymmetry
Let us now turn to a discussion of the asymmetry in the pion-correlated Fermi gas model. For its evaluation the axial-vector response function $`R_\text{T}^{}^{\text{V}A}`$ is needed. In a non-relativistic context the latter is related to the isovector component of the EM transverse response, $`R_{\text{A}V,I=1}^\text{T}`$, through the simple formula
$$R_\text{T}^{}^{\text{V}A}(q,\omega )=a_V\frac{G_A^{(1)}}{G_M^{(1)}}\frac{1}{\kappa }R_\text{T}^{I=1}(q,\omega ),$$
(157)
which can be extended into the relativistic regime via the prescription
$$R_\text{T}^{}^{\text{V}A}(q,\omega )a_V\frac{G_A^{(1)}}{G_M^{(1)}}\sqrt{\frac{\tau +1}{\tau }}R_\text{T}^{\text{I}=1}(q,\omega );$$
(158)
these agree to better than 2% in the momentum range 300 MeV/c $`<q<1`$ GeV/c for the RFG. Although not immediately apparent, Eq. (158) has been demonstrated to be preserved even in presence of pionic correlations. The effect of pionic correlations and MEC effects in this response is a “hardening” (a shift to higher $`\omega `$) of the peak of the response at intermediate values of $`q`$, which then fades away and even leads to a slight “softening” of the response at the highest momentum transfers considered.
We are now in the position to calculate the asymmetry $`𝒜`$, which is displayed in Fig. 21 as a function of $`\omega `$ for three values of $`q`$ and for forward ($`\theta =10^0`$) and backward ($`\theta =170^0`$) scattering angles.
Upon examining Fig. 21, we notice the significant effect occurring at moderate $`q`$ (say 300–500 MeV/c), small $`\omega `$ and forward angles. As previously discussed, it is related to the large negative value assumed by the correlated $`R_\text{L}^{\text{A}V}`$, which leads to a pionic asymmetry that is an order-of-magnitude larger than the free RFG one. However, as $`\omega `$ increases $`R_\text{L}^{\text{A}V}`$ rapidly decreases until it changes sign, while $`R_\text{T}^{\text{A}V}`$ stays negative; accordingly, they largely cancel in the numerator of the ratio expressing $`𝒜`$ and this becomes substantially lowered.
Interestingly, an energy is reached (about 60 MeV for $`q=300`$ MeV/c) where the correlated and free RFG values of $`𝒜`$ coincide. At still larger $`\omega `$ a further reduction of $`𝒜`$ is seen to occur until at about 90 MeV it nearly vanishes. This constitutes an example of a dynamical restoration of a symmetry (here the left-right parity symmetry) and reflects the complex nature of the PV longitudinal response. The near-vanishing of the asymmetry at $`\omega 90`$ MeV in Fig. 21 stems from the cancellation between the positive contribution it gets from $`R_\text{L}^{\text{A}V}`$ and the negative one it gets from $`R_\text{T}^{\text{A}V}`$. This trend of course fades away at larger $`\theta `$, where the role of the longitudinal PV response gradually becomes irrelevant. Finally, we see that at larger momenta, where the impact of correlations is no longer so strongly felt, the nearly perfect restoration of the left-right symmetry does not show up anymore. It is, however, still true (even at 1 GeV/c) that an energy exists where the free and the correlated values of the asymmetry coincide.
From the lower panels of Fig. 21 it clearly appears that at backward angles $`𝒜`$ is almost equally sensitive to the nuclear dynamical effects and to the uncertainties in the nucleonic form factors. This raises the question of whether or not the asymmetry itself is a suitable observable for extracting information on either the nuclear correlations or the single-nucleon form factors. As a consequence we are led to consider three different integrated observables that have been introduced to emphasize one of the two aspects of the problem. They are:
* $$\mathrm{\Delta }𝒜(q,\theta )\frac{1}{\mathrm{\Delta }\omega }\left[_{\omega _{min}}^{\omega _{QEP}}𝑑\omega 𝒜(\theta ;q,\omega )_{\omega _{QEP}}^{\omega _{max}}𝑑\omega 𝒜(\theta ;q,\omega )\right]$$
(159)
* $$\overline{𝒜}(q,\theta )\frac{1}{\mathrm{\Delta }\omega }_{\omega _{min}}^{\omega _{max}}𝑑\omega 𝒜(\theta ;q,\omega )$$
(160)
and
* $$(q,\theta )\frac{_{\omega _{min}}^{\omega _{max}}𝑑\omega W^{PV}(q,\omega )/\stackrel{~}{X}_\text{T}(\theta ,\tau ,\psi ;\eta _F)}{_{\omega _{min}}^{\omega _{max}}𝑑\omega W^{EM}(q,\omega )/X_\text{T}(\theta ,\tau ,\psi ;\eta _F)}.$$
(161)
Here $`\omega _{min}`$ and $`\omega _{max}`$ are the RFG response boundaries for a fixed $`q`$
$`\omega _{min}`$ $`=`$ $`\sqrt{(k_Fq)^2+m_N^2}\sqrt{k_F^2+m_N^2}`$ (162)
$`\omega _{max}`$ $`=`$ $`\sqrt{(k_F+q)^2+m_N^2}\sqrt{k_F^2+m_N^2},`$ (163)
the energy interval $`\mathrm{\Delta }\omega `$ is
$`\mathrm{\Delta }\omega `$ $``$ $`\omega _{max}\omega _{min}`$ (164)
$`=`$ $`\sqrt{(k_F+q)^2+m_N^2}\sqrt{(k_Fq)^2+m_N^2}`$ (165)
$``$ $`2k_Fq/\sqrt{q^2+m_N^2}`$ (166)
and the hadronic functions
$`W^{EM}`$ $`=`$ $`v_\text{L}R_\text{L}+v_\text{T}R_\text{T}`$ (167)
$`W^{PV}`$ $`=`$ $`v_\text{L}R_\text{L}^{\text{A}V}+v_\text{T}R_\text{T}^{\text{A}V}+v_\text{T}^{}R_\text{T}^{}^{\text{V}A}`$ (168)
are divided by
$`X_\text{T}(\theta ,\tau ,\psi ;\eta _F)`$ $`=`$ $`v_\text{T}\left(2\tau G_M^2+{\displaystyle \frac{G_E^2+\tau G_M^2}{1+\tau }}\mathrm{\Delta }\right)`$ (169)
$`\stackrel{~}{X}_\text{T}(\theta ,\tau ,\psi ;\eta _F)`$ $`=`$ $`a_Av_\text{T}\left(2\tau G_M\stackrel{~}{G}_M+{\displaystyle \frac{G_E\stackrel{~}{G}_E+\tau G_M\stackrel{~}{G}_M}{1+\tau }}\mathrm{\Delta }\right)`$ (170)
in order to extract the single-nucleon content from the many-body one . The quantity $`\mathrm{\Delta }`$ is given later in Eq. (193): since $`\mathrm{\Delta }\eta _F^21`$, the contributions containing $`\mathrm{\Delta }`$ may often safely be neglected. Removing the single-nucleon content in this way has the following advantages:
* The pionic correlations are particularly felt by $`\mathrm{\Delta }𝒜`$, as is clearly apparent from Fig. 22. Indeed, there we first observe that at $`q=300`$ MeV/c the results obtained with the free RFG model almost vanish because of the nearly perfect cancellation between the contributions where $`\omega <\omega _{QEP}`$ and those where $`\omega >\omega _{QEP}`$; however, at larger $`q`$ this cancellation becomes less complete, owing partly to the role played by the nucleonic form factors and partly to the RFG model itself, whose responses (in contrast to the non-relativistic case) become less and less symmetric as $`q`$ increases. It is also clear that the correlations, in particular the exchange diagram, dramatically alter the prediction of the free RFG, yielding a huge $`\mathrm{\Delta }𝒜_\pi `$ at small $`\theta `$. This result is simply interpreted by observing that of the nuclear responses that enter in the asymmetry the pion has its greatest effect on $`R_\text{L}^{\text{A}V}`$ and although in the RFG model the latter accounts only for at most about 10% of the total asymmetry (and this only in the forward direction), nevertheless the impact of the pionic correlations is violent enough to induce a large negative value of $`R_\text{L}^{\text{A}V}`$ at small $`\omega `$, which is in turn reflected in the large negative value of $`\mathrm{\Delta }𝒜_\pi `$ at small $`\theta `$ displayed in Fig. 22. We deduce from this that the characteristic behavior of $`\mathrm{\Delta }𝒜`$ with $`\theta `$ shown in the figure represents one of the most transparent signatures of pion-induced isoscalar correlations in nuclei (we recall that $`R_\text{T}^{}^{\text{V}A}`$ is purely isovector and that in $`R_\text{T}^{\text{A}V}`$ the isoscalar contribution is strongly suppressed — see Ref. ).
The observable $`\mathrm{\Delta }𝒜`$ has been specifically devised to enhance the signal for nuclear correlations and to minimize the sensitivity to the single-nucleon form factors. This can be inferred by observing the three (almost overlapping) lines for each family of curves in Fig. 22, corresponding to a variation of the strength of the effective axial-vector coupling, $`g_A^{(1)}`$, of $`\pm 10`$% around the canonical value $`g_A^{(1)}=1.26`$. The impact on $`\mathrm{\Delta }𝒜`$ of pionic correlations is seen to be more than an order-of-magnitude larger than that arising from variations in the axial-vector form factor. Similar results are found for the magnetic and electric strangeness form factors (see Ref. ).
* The energy-averaged asymmetry $`\overline{𝒜}`$, on the other hand, has been devised to minimize the sensitivity to the pionic correlations, as these tend to cancel out in the symmetrical integral in Eq. (160). The MEC contribution, on the other hand, does not average out in $`\overline{𝒜}`$, although in the $`ph`$ sector of the nuclear excitations it turns out to have a rather insignificant effect. Typical results are shown in Fig. 23 for $`q=`$ 500 MeV/c.
In particular, as seen in the expanded view in the figure, at very backward scattering angles where one might hope to determine the effective axial-vector coupling $`g_A^{(1)}`$ (again variations of $`\pm 10`$% around 1.26 are shown in the figure) the free RFG and pionic correlated results for the energy-averaged asymmetry come together (for $`q`$=500 MeV/c at $`\theta 147^o`$, which however varies with $`q`$). Since this special condition is presumably model-dependent, it is unlikely that one can count on using such particular kinematics to effect a determination of $`g_A^{(1)}`$ through the variations shown in the figure.
In contrast to the backward-angle situation, at forward scattering angles where the pionic correlations induce drastic modifications in $`R_\text{L}^{\text{A}V}`$, as we have seen, here the two families of curves differ although certainly not as much as in the case of $`\mathrm{\Delta }𝒜`$. In other words, the observable $`\overline{𝒜}`$ has some of the properties that we are looking for when we construct quantities that suppress the effects of correlations while bringing out the dependences on the single-nucleon form factors; however, this particular observable appears not to be entirely optimal. Since the PV longitudinal response is so strongly affected by the presence of pionic correlations, it is necessary to adopt an alternative approach to minimize these effects, leading to the introduction of the observable in Eq. (161).
* The procedure proposed in Refs. for the definition of the quantity $``$ has the goal of scaling the results (see the next section) by dividing out most of the single-nucleon content through the use of the dividing factors $`X_\text{T}`$ and $`\stackrel{~}{X}_\text{T}`$. In Fig. 24 we show $``$ as a function of $`\theta `$ for three values of $`q`$. Again two families of curves are displayed, one for the free RFG and one for the model with pionic effects included (labeled $`\pi `$), and each family has three curves ($`g_A^{(1)}=1.26`$, $`1.26+10`$% and $`1.2610`$%). In panel b we see a significant range of angles over which the pionic effects provide negligible modifications with respect to the RFG results and where the (merged) curves with the three values of the isovector/axial-vector strength can clearly be discerned. This behavior is similar at $`q=1`$ GeV/c, although not quite as nicely separated. Even so, the difference between the RFG and pionic families at backward scattering angles amounts to an effective change in $`g_A^{(1)}`$ of only about 4%.
It thus appears that the observable $``$ is sufficiently uncontaminated by correlation effects for favorable kinematics and better suited than $`\overline{𝒜}`$ to disentangling the nucleonic form factors from the nuclear dynamics. It has been specifically designed to integrate out the anomalous $`\omega `$-dependence, leaving quantities that for the most part only retain sensitivities to variations in the single-nucleon form factors.
## IV Scaling and sum rules
In this section we briefly address several issues that arise in discussing scaling and sum rules, showing how these two properties are interrelated and how they constrain the nuclear models. At sufficiently large three-momentum transfer $`q`$ the so-called $`y`$-scaling region occurs when the energy transfer $`\omega `$ is lower than its value at the quasielastic peak and when non-quasielastic processes such as meson production do not affect the nuclear responses. The $`y`$-scaling approach attempts to find some function of $`q`$ and $`\omega `$, here denoted $`G`$, such that when divided into the inclusive electron scattering cross section, the result is a reduced response $`F(q,y)`$ that scales as a function of $`y`$. Here $`y`$ is an appropriately chosen scaling variable (see below), is a function of $`(q,\omega )`$ and replaces $`\omega `$ (i.e., one uses the variables $`q`$ and $`y`$, rather than $`q`$ and $`\omega `$). “Scaling” means that for sufficiently large momentum transfers the function $`F`$ becomes universal, namely a function only of $`y`$, but not of $`q`$:
$$F(q,y)\stackrel{q\mathrm{}}{}F(y)F(\mathrm{},y).$$
(171)
The choice of the dividing function and scaling variable must be such as to remove the single-nucleon content from the nuclear responses in as model-independent a manner as possible while still retaining essential relativistic effects whenever feasible.
A parallel strategy concerns the Coulomb Sum Rule (CSR) : a dividing function $`H_\text{L}`$ can be devised such that the corresponding reduced longitudinal response $`r_\text{L}=R_\text{L}/H_\text{L}`$ fulfills the CSR.
In past work medium- and high-energy data have been tested with both the usual $`y`$-scaling approach (for a review, see Ref. ), while more recently the RFG-motivated approach has been applied and seen to scale successfully . Actually these data appear to support not only scaling, but also superscaling, namely the existence of a function related to $`F`$ that is the same for all nuclei . Additionally, it now appears that the experimental CSR is reasonably well saturated at high momentum transfers . As a consequence, any reliable nuclear model should simultaneously fulfill the two major requirements of
1) scaling
2) fulfilling the Coulomb Sum Rule.
We now show that the RFG model simultaneously satisfies the above properties. In the context of the PWIA for $`(e,e^{}N)`$ reactions the $`y`$-scaling variable is defined to be equal and opposite to the smallest value of the missing momentum, $`p_{min}`$, attained in the $`y`$-scaling region: It turns out to read
$`y`$ $`=`$ $`p_{min}={\displaystyle \frac{1}{2W^2}}\{(M_A^0+\omega )\sqrt{W^2\left(M_{A1}^0+m_N\right)^2}\sqrt{W^2\left(M_{A1}^0m_N\right)^2}`$ (173)
$`q[W^2+\left(M_{A1}^0\right)^2m_N^2]\},`$
where
$$W=\sqrt{\left(M_A^0+\omega \right)^2q^2},$$
(174)
$`M_A^0`$ and $`M_{A1}^0`$ being the masses of the initial and daughter nuclei (in their ground states), respectively. The energy transfer can, of course, be expressed in terms of $`q`$ and $`y`$. In particular, it must lie in the range $`\omega _t\omega q`$, where
$$\omega _t=E_S+\sqrt{(M_{A1}^0+m_N)^2+q^2}(M_{A1}^0+m_N)$$
(175)
is the threshold energy, with $`E_S=m_N+M_{A1}^0M_A^0`$ being the nuclear separation energy. The scaling variable in Eq. (173) vanishes when
$$\omega =\omega _0=E_S+\sqrt{m_N^2+q^2}m_N,$$
(176)
which is roughly the position of the quasielastic peak, and hence the scaling region is characterized by having $`y`$ negative.
In order to study the scaling behavior of the inclusive $`(e,e^{})`$ process one should be able to remove the effective $`eN`$ cross section from under the integrals involved in going from coincidence to inclusive scattering. These integrals extend over the missing momentum $`(p_m)`$ and over an energy that characterizes the degree of excitation of the daughter nucleus, $``$:
$$=E_{A1}E_{A1}^00,$$
(177)
where $`E_{A1}`$ is the energy of the unobserved daughter system (in general in an excited state) and $`E_{A1}^0`$ is that energy when this system is in its ground state, i.e., has mass $`M_{A1}^0`$. Naturally $``$ can be re-expressed in terms of the missing energy $`E_m`$; one has roughly that $`E_m+E_S`$. If one assumes that the proton and neutron distributions inside the nucleus are equal, which is a reasonable approximation for $`N=Z`$ nuclei, and that the most important contributions to the nuclear spectral function arise from the lowest values of $`(p,)`$ that can be reached for given values of $`q`$ and $`y`$ (in the scaling region these are $`=0`$ and $`p=y`$), then the function one hopes will scale as a function of $`y`$ when $`q\mathrm{}`$ is
$$F(q,y)\frac{d^2\sigma /d\mathrm{\Omega }_ed\omega }{\stackrel{~}{\sigma }_{eN}(q,y;p=y,=0)}.$$
(178)
In the PWIA this is indeed found to be the case, since in the limit $`q\mathrm{}`$ Eq. (178) becomes a function only of $`y`$, namely it scales .
Turning now to the RFG model, its spectral function is
$$\stackrel{~}{S}^{\text{RFG}}(p,)=\frac{3A}{8\pi k_F^3}\theta (k_Fp)\delta \left[(p)^{\text{RFG}}(p)\right],$$
(179)
where the excitation energy is
$$^{\text{RFG}}(p)=\left(\sqrt{k_F^2+m_N^2}\sqrt{p^2+m_N^2}\right).$$
(180)
Defining the RFG scaling variable through the intercept of the support of the RFG spectral function given in Eq. (179) and the kinematical boundaries in the missing energy-missing momentum plane one obtains the $`y`$-scaling variable of RFG
$$y_{\text{RFG}}=m_N\zeta =m_N\left(\lambda \sqrt{1+\frac{1}{\tau }}\kappa \right),$$
(181)
where the dimensionless variables $`\lambda =\omega /2m_N`$, $`\kappa =q/2m_N`$ and $`\tau =\kappa ^2\lambda ^2`$ have been introduced.
A different scaling variable was originally proposed for the RFG in Ref. , namely
$$\psi =\frac{1}{\sqrt{\xi _F}}\frac{\lambda \tau }{\sqrt{(1+\lambda )\tau +\kappa \sqrt{\tau (1+\tau )}}},$$
(182)
where $`\xi _F=ϵ_F1=\sqrt{1+\eta _F^2}1`$ and $`\eta _F=k_F/m_N`$ are the dimensionless Fermi kinetic energy and momentum, respectively. With some algebra it can be shown that the relation between the two scaling variables is
$$\xi _F\psi ^2=\sqrt{1+(y_{\text{R}FG}/m_N)^2}1.$$
(183)
The physical significance of $`\psi `$ is then immediately apparent: among the nucleons responding to an external probe one has the smallest kinetic energy and this is given by $`\psi ^2`$ (in units of the dimensionless Fermi kinetic energy $`\xi _F`$). Instead of working from the unseparated inclusive cross section towards a reduced response that, if successful, would scale as $`q\mathrm{}`$, one can work directly with the separated longitudinal and transverse responses, $`R_\text{L}`$ and $`R_\text{T}`$, since
* we are most interested in model-to-model comparisons and the same procedures may be followed in each case (i.e., focusing on L or T responses directly),
* a few cases exist where L/T separations have been performed experimentally,
* we wish to draw comparisons with studies of the CSR where only the $`L`$ response is relevant.
In this spirit we seek reduced responses denoted $`F_{\text{L},T}(\kappa ,\psi )`$ that scale. These are to be obtained from the inclusive response functions $`R_{\text{L},T}(\kappa ,\lambda )`$ by dividing through by specific functions, denoted $`G_{\text{L},T}(\kappa ,\lambda )`$:
$$F_{\text{L},T}(\kappa ,\psi )R_{\text{L},T}(\kappa ,\lambda )/G_{\text{L},T}(\kappa ,\lambda ).$$
(184)
If the dividing functions are chosen appropriately, then as above the reduced responses defined in Eq. (184) will scale, namely, become functions only of a single scaling variable such as $`\psi `$ defined above when $`\kappa \mathrm{}`$,
$$F_{\text{L},T}(\kappa ,\psi )\stackrel{\kappa \mathrm{}}{}F_{\text{L},T}(\psi )F_{\text{L},T}(\mathrm{},\psi ).$$
(185)
Such dividing functions, derived in Ref. , are
$`G_\text{L}(\kappa ,\lambda )`$ $`=`$ $`{\displaystyle \frac{ZU_{\text{L}p}+NU_{\text{L}n}}{2\kappa [1+\xi _F(1+\psi ^2)/2]}}`$ (186)
$`=`$ $`{\displaystyle \frac{1}{2\kappa }}(ZU_{\text{L}p}+NU_{\text{L}n})+𝒪(\xi _F)`$ (187)
and
$`G_\text{T}(\kappa ,\lambda )`$ $`=`$ $`{\displaystyle \frac{ZU_{\text{T}p}+NU_{\text{T}n}}{2\kappa [1+\xi _F(1+\psi ^2)/2]}}`$ (188)
$`=`$ $`{\displaystyle \frac{1}{2\kappa }}(ZU_{\text{T}p}+NU_{\text{T}n})+𝒪(\xi _F),`$ (189)
where (see Ref. )
$$U_{\text{L}p,n}=\frac{\kappa ^2}{\tau }\left[G_{\text{E}p,n}^2(\tau )+W_{2\text{p},n}(\tau )\mathrm{\Delta }\right]$$
(190)
$$U_{\text{T}p,n}=2\tau G_{\text{M}p,n}^2(\tau )+W_{2\text{p,n}}(\tau )\mathrm{\Delta },$$
(191)
with
$$W_{2\text{p},n}(\tau )=\frac{1}{1+\tau }\left[G_{\text{E}p,n}^2(\tau )+\tau G_{\text{M}p,n}^2(\tau )\right]$$
(192)
and
$$\mathrm{\Delta }=\frac{\tau }{\kappa ^2}\left[\frac{1}{3}\left(ϵ_F^2+ϵ_F\sqrt{1+\zeta ^2}+1+\zeta ^2\right)+\lambda \left(ϵ_F+\sqrt{1+\zeta ^2}\right)+\lambda ^2\right](1+\tau ).$$
(193)
Dividing the longitudinal and transverse RFG response functions by $`G_\text{L}`$ and $`G_\text{T}`$ yields the reduced responses
$$F_\text{L}^{\text{RFG}}(\psi )=F_\text{T}^{\text{RFG}}(\psi )=\frac{3\xi _F}{2m_N\eta _F^3}(1\psi ^2)\theta (1\psi ^2)\left[1+\frac{1}{2}\xi _F(1+\psi ^2)\right],$$
(194)
that, by construction, scale with $`\psi `$, $`\zeta `$ or $`y_{\text{R}FG}`$.
In parallel with the scaling behavior of the RFG one can study the CSR and the various energy-weighted moments of another reduced response denoted $`r_\text{L}(\kappa ,\lambda )`$, introduced in Refs. . Here the longitudinal response $`R_\text{L}(\kappa ,\lambda )`$ is divided by a function $`H_\text{L}(\kappa ,\lambda )`$ to yield
$$r_\text{L}(\kappa ,\psi )R_\text{L}(\kappa ,\lambda )/H_\text{L}(\kappa ,\lambda )$$
(195)
and the n<sup>th</sup> moment of the longitudinal response of the nucleus is given by
$$\mathrm{\Xi }^{(n)}=\underset{0}{\overset{\kappa }{}}𝑑\lambda \lambda ^nr_\text{L}(\kappa ,\lambda ).$$
(196)
In particular, the $`n=0`$ moment, $`\mathrm{\Xi }^{(0)}`$, is the CSR. In the case of the RFG the dividing function is
$$H_\text{L}(\kappa ,\lambda )=\frac{\kappa \eta _F^3}{2\xi _F}\left(ZU_{\text{L}p}+NU_{\text{L}n}\right)/\left(\psi /\lambda \right).$$
(197)
Thus, upon dividing the charge response of the RFG by Eq. (197) one obtains the following reduced longitudinal response
$$r_\text{L}^{\text{RFG}}(\kappa ,\lambda )=\frac{3}{8m_N}(1\psi ^2)\theta (1\psi ^2)\frac{\psi }{\lambda },$$
(198)
which, by construction, fulfills the CSR in the non-Pauli-blocked domain, as can easily be verified.
The two dividing functions $`G_\text{L}`$ (related to the scaling) and $`H_\text{L}`$ (related to the CSR) are linked according to
$`G_\text{L}(\kappa ,\lambda )`$ $`=`$ $`\left({\displaystyle \frac{\eta _F^3}{4\xi _F}}\right){\displaystyle \frac{\psi }{\lambda }}{\displaystyle \frac{1}{1+\xi _F(1+\psi ^2)/2}}H_\text{L}(\kappa ,\lambda )`$ (199)
$`=`$ $`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\kappa }{\tau }}\right)\left({\displaystyle \frac{1+2\lambda }{1+\lambda }}\right)H_\text{L}(\kappa ,\lambda )+𝒪(\xi _F).`$ (200)
When experimental data are reduced using the above dividing functions, they are seen to yield a CSR at high-$`q`$ and to scale when plotted versus any of the scaling variables previously introduced.
Finally we explore the scaling and CSR properties of nuclear models other than the RFG, specifically:
* the hybrid model (HM), introduced in Ref. ;
* the quantum hadrodynamical (QHD) model .
The hybrid model is designed to account for the binding of the nucleons inside the nucleus, thus curing a flaw of the RFG related to its negative separation energy. The HM has continuum states that are plane waves, as in the RFG model, but has bound states described by shell-model wave functions obtained by solving the Schrödinger equation with some choice of potential well (in harmonic oscillator bound-state wave functions were used to simplify the analysis in the limit where $`A\mathrm{}`$). In the HM the scaling variable cannot be obtained analytically, since the model can only be treated numerically. However, scaling variables exist that incorporate the shift result from using Eq. (181) for $`\zeta ^{}`$ (and a corresponding dimensionful variable $`y^{}`$) or Eq. (182) for $`\psi ^{}`$ by making the replacements $`\lambda \lambda ^{}`$ and $`\tau \tau ^{}=\kappa ^2\lambda _{}^{}{}_{}{}^{2}`$, where
$$\lambda ^{}=\lambda \lambda _{\text{shift}},$$
(201)
with
$$\lambda _{\text{shift}}=\frac{1}{2m_N}(T_F+E_S)$$
(202)
and $`T_F=m_N\xi _F`$ the Fermi kinetic energy. The HM turns out to have the width of its reduced response identical to a RFG computed with a Fermi momentum that is somewhat larger than the usual one (237 MeV/c for the HM, versus the 230 MeV/c value for $`k_F`$ used for the RFG to correspond to nuclei near <sup>40</sup>Ca or <sup>56</sup>Fe).
Scaling may then be examined for the HM by computing
$$F_\text{L}^{\text{HM}}(\kappa ,\psi )R_\text{L}^{\text{HM}}(\kappa ,\lambda )/G_L(\kappa ,\lambda )$$
(203)
and the various energy-weighted moments of the longitudinal response, including the zeroth moment or CSR, may be computed using
$$r_\text{L}^{\text{HM}}(\kappa ,\psi )R_\text{L}^{\text{HM}}(\kappa ,\lambda )/H_\text{L}(\kappa ,\lambda ).$$
(204)
Note that the same dividing factors $`G_\text{L}`$ and $`H_\text{L}`$ that were developed from our discussions of the RFG are used.
In Fig. 25 the scaling function $`F_\text{L}^{\text{HM}}`$ for <sup>40</sup>Ca is displayed versus $`\psi `$ and $`\psi ^{}`$ for four different values of $`q`$ (the RFG result is also shown for reference): the HM scales either with $`\psi `$ or with $`\psi ^{}`$ as $`q`$ becomes large. Indeed, only the $`q=`$ 500 MeV/c plot versus $`\psi `$ shows any appreciable violation of scaling, whereas the scaling versus $`\psi ^{}`$ is excellent.
In the QHD model protons and neutrons in the nucleus are described by Dirac spinors and move in strong Lorentz scalar and vector mean fields. These in turn arise self-consistently from the exchange of $`\sigma `$ and $`\omega `$ mesons between the same nucleons on which they act. The scalar field dresses the bare mass of the nucleon, considerably lowering its value; the vector field uniformly shifts the fermion spectrum. As a consequence the QHD charge response of nuclear matter in Hartree approximation is unaffected by the vector field, while it turns out to be quite sensitive to the effective mass $`m_N^{}`$ induced by the scalar field. This is, of course, true in the simple approximation of constant relativistic mean fields. An improved description allows for an energy-dependence of the latter, which helps to account for the data of proton-nucleus elastic scattering.
As for the HM model, scaling can be examined in the QHD model by computing
$$F_\text{L}^{\text{QHD}}(\kappa ,\psi )R_\text{L}^{\text{QHD}}(\kappa ,\lambda )/G_\text{L}(\kappa ,\lambda )$$
(205)
and likewise the various energy-weighted moments of the longitudinal response computed using
$$r_\text{L}^{\text{QHD}}(\kappa ,\psi )R_\text{L}^{\text{QHD}}(\kappa ,\lambda )/H_\text{L}(\kappa ,\lambda ).$$
(206)
Since the dividing factors are (at least to a very good level of approximation) universal, accordingly we use the same dividing factors $`G_\text{L}`$ and $`H_\text{L}`$ that were developed from our discussions of the RFG.
In Fig. 26 we display the reduced responses in Eqs. (205) and (206) as functions both of $`\psi `$ and also $`\psi ^{}`$, namely, the RFG scaling variable given in Eq. (182) with $`m_N`$ replaced by $`m_N^{}`$ (two different values of the effective mass, $`m_N^{}=0.68m_N`$ and 0.8 $`m_N`$ are used). It is clearly seen that $`F_\text{L}^{\text{QHD}}`$ does not scale versus $`\psi `$ when the effective mass is constant and differs from $`m_N`$. As $`q`$ continues to grow beyond the range of values shown in the figures, the results continue to shift to higher $`\omega `$ and never coalesce into a universal curve. When plotted versus $`\psi ^{}`$ the behavior, while better, still does not scale. This is in contrast with the RFG and HM results displayed above and, importantly, is not what is seen experimentally where the world data do appear to scale in $`\psi `$ . The fact that experimentally the scaling is observed to occur successfully for $`q`$ greater than about 1 GeV/c suggests that $`m_N^{}/m_N`$ should not deviate appreciably from unity for such kinematics.
Finally, the CSR ($`\mathrm{\Xi }^{(0)}`$), the energy-weighted sum rule ($`\mathrm{\Xi }^{(1)}/\mathrm{\Xi }^{(0)}`$) and the variance $`\sigma =\sqrt{\mathrm{\Xi }^{(2)}(\mathrm{\Xi }^{(1)})^2}`$ of <sup>40</sup>Ca corresponding to the HM, QHD (with $`m_N^{}=0.68`$ $`m_N`$ and 0.8 $`m_N`$) and RFG models are displayed in Fig. 27. The RFG model and the HM both saturate the CSR at high-$`q`$. In contrast, the QHD model does so only if the effective value of $`m_N^{}/m_N`$ evolves with increasing $`q`$ towards unity, as suggested by the latest version of the model.
## V Outlook and perspectives
In this paper we have discussed the parity-conserving and violating- inclusive nuclear responses, addressing the following issues
* Where do they occur?
* How can we describe them?
* What is their input?
With respect to the first item we have limited our focus to the quasielastic peak (QEP), since at larger excitation energies where for example the $`\mathrm{\Delta }`$ plays a role we do not expect a description of the nucleus only in terms of nucleonic and mesonic degrees of freedom to be adequate. At some point QCD degrees of freedom should become the more appropriate ones to describe inclusive scattering. On the other hand at energies significantly lower than those characterizing the QEP where discrete excitations, giant resonances, etc. are seen the theoretical many-body framework is different from the one used here.
Concerning the second item let us again stress that any theoretical framework should first fulfill (as much as possible) Lorentz covariance and gauge invariance. For this the RFG appears to be a good starting point: It is a covariant model, since its ingredients are the fully relativistic nucleon propagators and EM (or WNC) vertices, and it respects gauge invariance, since the vector currents of the nucleon are conserved. When mesons are added to the picture then Lorentz covariance and gauge invariance are fulfilled to the extent that one allows for a dynamical propagation of the mesons and treats the forces and the currents consistently. This turns out to be possible for the case of the pion.
Of course the RFG misses surface and finite-size effects; however, first of all, these are of minor relevance for the scattering of electrons in the QEP and $`\mathrm{\Delta }`$ peak domains, and secondly, they can be satisfactorily accounted for within the semiclassical approach, which exploits the advantages offered by the translational invariance of the RFG and accommodates these advantages to fit the physics of finite systems .
In the framework of the RFG we have treated nucleon-nucleon correlations in a perturbative scheme. Alternatives are represented by
* variational approaches
* loop expansions in the path integral framework, which however represent an alternative regrouping of perturbation theory.
Concerning perturbation theory we have seen that basic ingredients are the diagrams associated with the HF mean field and those related to the antisymmetrized RPA. The emphasis on the antisymmetrization arises from the recognition that very important carriers of the nuclear force, notably $`\pi `$ and $`\rho `$ mesons, only act through the exchange diagrams.
The short-range correlations (SRC) induced by the violent repulsion present in the N-N force at small distances are then inserted via the ladder diagrams. It is, however, doubtful whether ladders can be covariantly computed, especially at high densities (or Fermi momenta $`k_F`$) where they become increasingly important and the role of relativity cannot be ignored. In addition, it may be that searching for a totally deterministic account of the N-N correlations is not the optimal way to proceed — should, for example, quantum chaos be at work in atomic nuclei, then costly efforts to compute SRC would not represent the most efficient way of interpreting the nuclear response functions.
A further serious challenge facing the perturbative approach relates to the estimate of the size of diagrams that are not considered. In this connection it appears that only the loop expansion offers a consistent criterion (the number of loops) to organize the perturbative series in classes of homogeneous diagrams allowing at the same time an estimate of the rate of convergence of the resulting new expansion. However, although this is theoretically established, its practical implementation is far from trivial.
Finally a few words on the input to be fed into the perturbative scheme are appropriate. In general the Bonn potential appears as a well-founded representation of the N-N force. Indeed it provides an excellent representation of the energy behavior of the N-N phase shifts in free space and of the deuteron’s properties. In addition, it can be derived from a Lagrangian defined in terms of nucleonic and mesonic fields, and therefore cast in the framework of an effective field theory that, from the theoretical viewpoint, represents a highly desirable feature. In particular, it allows for a test of the gauge invariance of the theory, i.e. the consistency between the forces and currents. In this connection it is worthwhile to emphasize that the most direct way to assess the validity of the mesons-plus-nucleons model of the nucleus ultimately rests on a deeper understanding of the role of meson exchange currents (MEC) in the nuclear responses. Much has been done in this field, although much remains to be done.
It seems to us that at present the analysis of the nuclear responses is best performed on the basis of a hadronic approach such as the one we have pursued using the Bonn potential. It is important to remember, however, that such potentials are being used for nucleons within the nucleus and that these are differently off-shell from the conditions found in studying N-N scattering or the ground state of the deuteron.
In summary, in Section II of the present article we have analyzed the quasielastic response functions for inclusive electron scattering in the standard framework of HF plus RPA. In addition, the impact of short-range correlations on the effective particle-hole force has been explored within the context of the G-matrix approach. The analysis has then been extended in Section III to include the observables that occur in studying parity-violating electron scattering, with emphasis placed on the “new” nuclear axial response. In Section IV the attention has been directed towards several general, important properties of the response functions, in particular their Coulomb sum rule and scaling/superscaling behavior. Finally, a short account of what lies ahead has been provided in Section V.
## A Electromagnetic form factors of the nucleon
Here we give the formulae for the EM form factors introduced in the definition of the response functions in Eq. (19).
### 1 Non-relativistic Fermi gas
In a non-relativistic calculation one defines, in terms of Sach’s form factors,
$`f_{L}^{(I)}{}_{}{}^{2}`$ $`=`$ $`G_{E}^{(I)}{}_{}{}^{2}`$ (A2)
$`f_{T}^{(I)}{}_{}{}^{2}`$ $`=`$ $`2\tau G_{M}^{(I)}{}_{}{}^{2},I=0,1,`$ (A3)
where $`\tau =|Q^2|/4m_N^2=(q^2\omega ^2)/4m_N^2`$, $`G_X^{(I)}=G_{X_p}+(1)^IG_{X_n}`$ ($`X=E,M`$). A typical parameterization (although many others are possible) is the dipolar-plus-Galster one, namely
$$\begin{array}{ccc}G_{E_p}(\tau )\hfill & =& G_D^V(\tau )\hfill \\ G_{M_p}(\tau )\hfill & =& \mu _pG_D^V(\tau )\hfill \\ G_{M_n}(\tau )\hfill & =& \mu _nG_D^V(\tau )\hfill \\ G_{E_n}(\tau )\hfill & =& \mu _n\tau G_D^V(\tau )\xi _n(\tau ).\hfill \end{array}$$
(A4)
Here $`G_D^V(\tau )=(1+\lambda _D^V\tau )^2`$ is the vector dipole form factor, with $`\lambda _D^V4.97`$, whereas $`\mu _p2.793`$ and $`\mu _n1.913`$ are the proton and neutron magnetic moments, respectively. For $`G_{E_n}`$ we have adopted the Galster parameterization with $`\xi _n(\tau )=(1+\lambda _n\tau )^1`$ and $`\lambda _n5.6`$.
### 2 Relativistic Fermi gas
In Ref. it was shown that in an RFG calculation a very good approximation to the exact treatment of the EM vertices can be obtained with the following definitions:
$`f_{\text{L}}^{(I)}{}_{}{}^{2}`$ $`=`$ $`\left[{\displaystyle \frac{1}{1+\tau }}G_{E}^{(I)}{}_{}{}^{2}+\tau G_{M}^{(I)}{}_{}{}^{2}{\displaystyle \frac{k_F^2}{2m_N^2}}(1\psi _r)^2\right],`$ (A6)
$`f_{\text{T}}^{(I)}{}_{}{}^{2}`$ $`=`$ $`{\displaystyle \frac{2\tau }{(1+\omega /2m_N)^2}}G_{M}^{(I)}{}_{}{}^{2},`$ (A7)
where $`\psi _r`$ is the relativistic scaling variable in Eq. (24), $`k_F`$ is the Fermi momentum and the other quantities have been defined in Eqs. (A4).
## B First-order self-energy
Here we give the analytic expressions for the first-order self-energy based on the potential in Eqs. (2)–(*II A). $`\mathrm{\Sigma }^{(1)}(k)`$ is the sum of direct (“Hartree”) and exchange (“Fock”) terms, namely
$$\mathrm{\Sigma }^{(1)}(k)\mathrm{\Sigma }^\text{H}(k)+\mathrm{\Sigma }^\text{F}(k),$$
(B1)
where
$$\mathrm{\Sigma }^\text{H}(k)=\rho V_0(0)$$
(B2)
and
$$\mathrm{\Sigma }^\text{F}(k)=\frac{3}{8}\rho \underset{\alpha }{}C_\text{F}^{(\alpha )}𝒮_\alpha ^\text{F}(k),$$
(B3)
$`\rho =2k_F^3/3\pi ^2`$ being the nuclear density. In the last equation we have introduced the spin-isospin coefficients (note that the tensor channels do not contribute)
$`C_\text{F}^{(0)}`$ $`=`$ $`1,C_\text{F}^{(\tau )}=3,C_\text{F}^{(\sigma )}=3,C_\text{F}^{(\sigma \tau )}=9,`$ (B4)
$`C_\text{F}^{(t)}`$ $`=`$ $`C_\text{F}^{(t\tau )}=0,`$ (B5)
and defined
$$𝒮_\alpha ^\text{F}(k)=\frac{1}{2\pi }𝑑𝒌^{}\theta (k_Fk^{})V_\alpha (𝒌𝒌^{}).$$
(B6)
In any non-tensor channel $`\alpha `$ the potential is expressed as a combination of the “$`\delta `$” and “momentum-dependent” pieces in Eq. (*II A), for which one finds
$$𝒮_\delta ^\text{F}(k)=\{\begin{array}{cc}g_\delta \frac{2}{3},\hfill & \mathrm{}=0\hfill \\ g_\delta (\lambda ^2\mu ^2)w_a^\text{F}(\lambda |k),\hfill & \mathrm{}=1\hfill \\ g_\delta (\lambda ^2\mu ^2)^2w_b^\text{F}(\lambda |k),\hfill & \mathrm{}=2\hfill \end{array}$$
(B7)
and
$$𝒮_{\text{MD}}^\text{F}(k)=\{\begin{array}{cc}g_{\text{MD}}\mu ^2w_a^\text{F}(\lambda |k),\hfill & \mathrm{}=0\hfill \\ g_{\text{MD}}\mu ^2[w_a^\text{F}(\mu |k)w_a^\text{F}(\lambda |k)],\hfill & \mathrm{}=1\hfill \\ g_{\text{MD}}\mu ^2[w_a^\text{F}(\mu |k)w_a^\text{F}(\lambda |k)(\lambda ^2\mu ^2)w_b^\text{F}(\lambda |k)],\hfill & \mathrm{}=2,\hfill \end{array}$$
(B8)
where $`\mathrm{}`$ represents the power of the form factors (see Eq. (*II A)). Here we have introduced the dimensionless form factor cut-off, $`\lambda =\mathrm{\Lambda }/k_F`$, and meson mass, $`\mu =m/k_F`$, and we have defined
$`w_a^\text{F}(\lambda |k)`$ $`=`$ $`1\lambda \left[\mathrm{arctan}\left({\displaystyle \frac{1k}{\lambda }}\right)+\mathrm{arctan}\left({\displaystyle \frac{1+k}{\lambda }}\right)\right]{\displaystyle \frac{\lambda ^2k^2+1}{4k}}\mathrm{ln}\left|{\displaystyle \frac{\lambda ^2+(k1)^2}{\lambda ^2+(k+1)^2}}\right|,`$ (B10)
$`w_b^\text{F}(\lambda |k)`$ $`=`$ $`{\displaystyle \frac{1}{2\lambda }}\left[\mathrm{arctan}\left({\displaystyle \frac{1k}{\lambda }}\right)+\mathrm{arctan}\left({\displaystyle \frac{1+k}{\lambda }}\right)\right]+{\displaystyle \frac{1}{4k}}\mathrm{ln}\left|{\displaystyle \frac{\lambda ^2+(k1)^2}{\lambda ^2+(k+1)^2}}\right|.`$ (B12)
## C Tensor interaction in the exchange diagrams
The $`n`$-th order exchange polarization propagator in presence of tensor interactions has an expression that is slightly more complicated than that in Eq. (110), because the tensor operators do not in general allow for a factorization of the azimuthal integrations. A generic diagram with $`m`$ non-tensor and $`nm`$ tensor interaction lines can instead be written as
$`\mathrm{\Pi }_{\alpha _1\mathrm{}\alpha _m,\alpha _{m+1}\mathrm{}\alpha _n}^{(n)\text{ex}}(q,\omega )`$ $`=`$ $`(1)^n\left({\displaystyle \frac{m_N}{q}}\right)^{n+1}\left({\displaystyle \frac{k_F}{2\pi }}\right)^{2n+2}`$ (C5)
$`\times {\displaystyle _1^1}dy_1{\displaystyle \frac{1}{2}}{\displaystyle _0^{1y_1^2}}dx_1\mathrm{}{\displaystyle _1^1}dy_{n+1}{\displaystyle \frac{1}{2}}{\displaystyle _0^{1y_{n+1}^2}}dx_{n+1}`$
$`\times {\displaystyle \frac{1}{\psi y_1+i\eta _\omega }}W_{\alpha _1}(x_1,y_1;x_2,y_2)\mathrm{}W_{\alpha _m}(x_m,y_m;x_{m+1},y_{m+1})`$
$`\times W_{\alpha _{m+1}\mathrm{}\alpha _n}(x_{m+1},y_{m+1};\mathrm{};x_{n+1},y_{n+1}){\displaystyle \frac{1}{\psi y_{n+1}+i\eta _\omega }}`$
$`+{\displaystyle (\omega \omega )},`$
where $`W_{\alpha _i}`$ has been defined for the non-tensor channels in Eq. (107) and
$`W_{\alpha _{m+1}\mathrm{}\alpha _n}(x_{m+1},y_{m+1};\mathrm{};x_{n+1},y_{n+1})=2^{nm}{\displaystyle \underset{ij}{}}{\displaystyle \underset{l_1\mathrm{}l_{nm}}{}}\mathrm{\Lambda }_{ji}{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\phi _{m+1}}{2\pi }}\mathrm{}{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\phi _{n+1}}{2\pi }}`$ (C6)
$`\times V_{\alpha _{m+1}}(𝒌_{m+1}𝒌_{m+2})S_{il_1}(\widehat{𝒌_{m+1}𝒌_{m+2}})\mathrm{}V_{\alpha _n}(𝒌_n𝒌_{n+1})S_{l_{nm}j}(\widehat{𝒌_n𝒌_{n+1}}).`$ (C7)
In the last expression we have introduced the tensors
$$S_{ij}(\widehat{𝒌})=3\widehat{𝒌}_i\widehat{𝒌}j\delta _{ij},$$
(C8)
such that $`_{ij}\sigma _i\sigma _jS_{ij}(\widehat{𝒌})=S_{12}(\widehat{𝒌})`$.
The first-order case is rather simple, since one again gets Eqs. (LABEL:eq:Pi1ex)–(117) with
$$W_\alpha (x,y;x^{},y^{})=_0^{2\pi }\frac{d\phi }{2\pi }V_\alpha (𝒌𝒌^{})S_{zz}(\widehat{𝒌𝒌^{}}).$$
(C9)
At second order, however, one can use Eqs. (123)–(124) only when just one tensor interaction is present.
## D First- and second-order exchange diagrams
Here we give the explicit expressions for the first- and second-order exchange diagrams, based on the potential in Eqs. (2)–(*II A). In Eqs. (LABEL:eq:Pi1ex) and (116) we have seen that
$$\mathrm{\Pi }_\alpha ^{(1)\text{ex}}(q,\omega )=\left(\frac{m_N}{q}\right)^2\frac{k_F^4}{(2\pi )^4}\left[𝒬_\alpha ^{(1)}(0,\psi )𝒬_\alpha ^{(1)}(\overline{q},\psi )+𝒬_\alpha ^{(1)}(0,\psi +\overline{q})𝒬_\alpha ^{(1)}(\overline{q},\psi +\overline{q})\right],$$
(D1)
where
$$𝒬_\alpha ^{(1)}(\overline{q},\psi )=2_1^1𝑑y\frac{1}{\psi y+i\eta _\omega }_1^1𝑑y^{}W_{\alpha }^{}{}_{}{}^{\prime \prime }(y,y^{};\overline{q})\frac{1}{yy^{}+\overline{q}},$$
(D2)
whereas from Eqs. (123) and (124) one has
$`\mathrm{\Pi }_{\alpha \alpha ^{}}^{(2)\text{ex}}(q,\omega )=\left({\displaystyle \frac{m_N}{q}}\right)^3{\displaystyle \frac{k_F^6}{(2\pi )^6}}[𝒬_{\alpha \alpha ^{}}^{(2)}(0,0;\psi )𝒬_{\alpha \alpha ^{}}^{(2)}(0,\overline{q};\psi )𝒬_{\alpha \alpha ^{}}^{(2)}(\overline{q},0;\psi )+𝒬_{\alpha \alpha ^{}}^{(2)}(\overline{q},\overline{q};\psi )`$ (D3)
$`𝒬_{\alpha \alpha ^{}}^{(2)}(0,0;\psi +\overline{q})+𝒬_{\alpha \alpha ^{}}^{(2)}(0,\overline{q};\psi +\overline{q})+𝒬_{\alpha \alpha ^{}}^{(2)}(\overline{q},0;\psi +\overline{q})𝒬_{\alpha \alpha ^{}}^{(2)}(\overline{q},\overline{q};\psi +\overline{q})],`$ (D4)
where
$$𝒬_{\alpha \alpha ^{}}^{(2)}(\overline{q}_1,\overline{q}_2;\psi )=_1^1𝑑y\frac{1}{2}_0^{1y^2}𝑑x𝒢_\alpha (x,y+\overline{q}_1;\psi +\overline{q}_1)\frac{1}{\psi y+i\eta _\omega }𝒢_\alpha ^{}(x,y+\overline{q}_2;\psi +\overline{q}_2)$$
(D6)
and
$$𝒢_\alpha (x,y;\psi )=_1^1𝑑y^{}\frac{1}{\psi y^{}+i\eta _\omega }W_\alpha ^{}(x,y;y^{}).$$
(D7)
For a meson-exchange potential the quantities that can be calculated analytically are those given by Eqs. (107), (C9), (127) and (117), namely
$`W_\alpha (x,y;x^{},y^{})`$ $`=`$ $`{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\phi }{2\pi }}V_\alpha (𝒌𝒌^{})\text{(non-tensor)}`$ (D9)
$`W_\alpha (x,y;x^{},y^{})`$ $`=`$ $`{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\phi }{2\pi }}V_\alpha (𝒌𝒌^{})S_{zz}(\widehat{𝒌𝒌^{}})\text{(tensor)}`$ (D10)
and
$`W_\alpha ^{}(x,y;y^{})`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _0^{1y_{}^{}{}_{}{}^{2}}}𝑑x^{}W_\alpha (x,y;x^{},y^{})`$ (D11)
$`W_\alpha ^{\prime \prime }(y,y^{};\overline{q})`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _0^{1y^2}}𝑑x{\displaystyle \frac{1}{2}}{\displaystyle _0^{1y_{}^{}{}_{}{}^{2}}}𝑑x^{}W_\alpha (x,y+\overline{q};x^{},y^{}).`$ (D12)
In any channel $`\alpha `$ the potential is expressed as a combination of the terms displayed in Eq. (*II A). Then, for each of them one finds
$`W_\delta (x,y;x^{},y^{})`$ $`=`$ $`\{\begin{array}{cc}g_\delta ,\hfill & \mathrm{}=0\hfill \\ g_\delta (\lambda ^2\mu ^2)w_a(\lambda |x,y;x^{},y^{}),\hfill & \mathrm{}=1\hfill \\ g_\delta (\lambda ^2\mu ^2)^2w_b(\lambda |x,y;x^{},y^{}),\hfill & \mathrm{}=2\hfill \end{array}`$ (D17)
$`W_{\text{MD}}(x,y;x^{},y^{})`$ $`=`$ $`\{\begin{array}{cc}g_{\text{MD}}\mu ^2w_a(\mu |x,y;x^{},y^{}),\hfill & \mathrm{}=0\hfill \\ g_{\text{MD}}\mu ^2[w_a(\mu |x,y;x^{},y^{})w_a(\lambda |x,y;x^{},y^{})],\hfill & \mathrm{}=1\hfill \\ g_{\text{MD}}\mu ^2[w_a(\mu |x,y;x^{},y^{})w_a(\lambda |x,y;x^{},y^{})\hfill & \\ (\lambda ^2\mu ^2)w_b(\lambda |x,y;x^{},y^{})],\hfill & \mathrm{}=2\hfill \end{array}`$ (D22)
$`W_{\text{TN}}(x,y;x^{},y^{})`$ $`=`$ $`\{\begin{array}{cc}g_{\text{TN}}\{[3(yy^{})^2+\mu ^2]w_a(\mu |x,y;x^{},y^{})1\},\hfill & \mathrm{}=0\hfill \\ g_{\text{TN}}\{[3(yy^{})^2+\mu ^2]w_a(\mu |x,y;x^{},y^{})\hfill & \\ [3(yy^{})^2+\lambda ^2]w_a(\lambda |x,y;x^{},y^{})\},\hfill & \mathrm{}=1\hfill \\ g_{\text{TN}}\{[3(yy^{})^2+\mu ^2]\hfill & \\ \times [w_a(\mu |x,y;x^{},y^{})w_a(\lambda |x,y;x^{},y^{})]\hfill & \\ (\lambda ^2\mu ^2)[3(yy^{})^2+\lambda ^2]w_b(\lambda |x,y;x^{},y^{})\},\hfill & \mathrm{}=2,\hfill \end{array}`$ (D29)
where again $`\mathrm{}`$ labels the power of the form factors, we have introduced the dimensionless form factor cut-off, $`\lambda =\mathrm{\Lambda }/k_F`$, and meson mass, $`\mu =m/k_F`$, and we have defined
$`w_a(\lambda |x,y;x^{},y^{})`$ $`=`$ $`\{[\lambda ^2+(yy^{})^2+x+x^{}]^24xx^{}\}^{1/2}`$ (D31)
$`w_b(\lambda |x,y;x^{},y^{})`$ $`=`$ $`{\displaystyle \frac{\lambda ^2+(yy^{})^2+x+x^{}}{\{[\lambda ^2+(yy^{})^2+x+x^{}]^24xx^{}\}^{3/2}}}.`$ (D32)
For $`W_\alpha ^{}`$ one finds
$`W_\delta ^{}(x,y;y^{})`$ $`=`$ $`\{\begin{array}{cc}g_\delta (1y_{}^{}{}_{}{}^{2})/2,\hfill & \mathrm{}=0\hfill \\ g_\delta (\lambda ^2\mu ^2)w_a^{}(\lambda |x,y;y^{}),\hfill & \mathrm{}=1\hfill \\ g_\delta (\lambda ^2\mu ^2)^2w_b^{}(\lambda |x,y;y^{}),\hfill & \mathrm{}=2\hfill \end{array}`$ (D37)
$`W_{\text{MD}}^{}(x,y;y^{})`$ $`=`$ $`\{\begin{array}{cc}g_{\text{MD}}\mu ^2w_a^{}(\mu |x,y;y^{}),\hfill & \mathrm{}=0\hfill \\ g_{\text{MD}}\mu ^2[w_a^{}(\mu |x,y;y^{})w_a^{}(\lambda |x,y;y^{})],\hfill & \mathrm{}=1\hfill \\ g_{\text{MD}}\mu ^2[w_a^{}(\mu |x,y;y^{})w_a^{}(\lambda |x,y;y^{})\hfill & \\ (\lambda ^2\mu ^2)w_b^{}(\lambda |x,y;y^{})],\hfill & \mathrm{}=2\hfill \end{array}`$ (D42)
$`W_{\text{TN}}^{}(x,y;y^{})`$ $`=`$ $`\{\begin{array}{cc}g_{\text{TN}}\{[3(yy^{})^2+\mu ^2]w_a^{}(\mu |x,y;y^{})(1y_{}^{}{}_{}{}^{2})/2\},\hfill & \mathrm{}=0\hfill \\ g_{\text{TN}}\{[3(yy^{})^2+\mu ^2]w_a^{}(\mu |x,y;y^{})\hfill & \\ [3(yy^{})^2+\lambda ^2]w_a^{}(\lambda |x,y;y^{})\},\hfill & \mathrm{}=1\hfill \\ g_{\text{TN}}\{[3(yy^{})^2+\mu ^2][w_a^{}(\mu |x,y;y^{})w_a^{}(\lambda |x,y;y^{})]\hfill & \\ (\lambda ^2\mu ^2)[3(yy^{})^2+\lambda ^2]w_b^{}(\lambda |x,y;y^{})\},\hfill & \mathrm{}=2,\hfill \end{array}`$ (D48)
where
$`w_a^{}(\lambda |x,y;y^{})={\displaystyle \frac{1}{2}}`$ (D50)
$`\times \mathrm{ln}\left|{\displaystyle \frac{\lambda ^2+(yy^{})^2+1y_{}^{}{}_{}{}^{2}x+\sqrt{[\lambda ^2+(yy^{})^2+1y_{}^{}{}_{}{}^{2}+x]^24(1y_{}^{}{}_{}{}^{2})x}}{2[\lambda ^2+(yy^{})^2]}}\right|`$ (D51)
(D52)
$`w_b^{}(\lambda |x,y;y^{})={\displaystyle \frac{1}{4}}{\displaystyle \frac{1}{\lambda ^2+(yy^{})^2}}\left[1{\displaystyle \frac{\lambda ^2+(yy^{})^21+y_{}^{}{}_{}{}^{2}+x}{\sqrt{[\lambda ^2+(yy^{})^2+1y_{}^{}{}_{}{}^{2}+x]^24(1y_{}^{}{}_{}{}^{2})x}}}\right].`$ (D53)
Finally, for $`W_\alpha ^{\prime \prime }`$ one finds
$`W_\delta ^{\prime \prime }(y,y^{};\overline{q})`$ $`=`$ $`\{\begin{array}{cc}g_\delta [(1y^2)/2][(1y_{}^{}{}_{}{}^{2})/2],\hfill & \mathrm{}=0\hfill \\ g_\delta (\lambda ^2\mu ^2)w_a^{\prime \prime }(\lambda |y,y^{};\overline{q}),\hfill & \mathrm{}=1\hfill \\ g_\delta (\lambda ^2\mu ^2)^2w_b^{\prime \prime }(\lambda |y,y^{};\overline{q}),\hfill & \mathrm{}=2\hfill \end{array}`$ (D59)
$`W_{\text{MD}}^{\prime \prime }(y,y^{};\overline{q})`$ $`=`$ $`\{\begin{array}{cc}g_{\text{MD}}\mu ^2w_a^{\prime \prime }(\mu |y,y^{};\overline{q}),\hfill & \mathrm{}=0\hfill \\ g_{\text{MD}}\mu ^2[w_a^{\prime \prime }(\mu |y,y^{};\overline{q})w_a^{\prime \prime }(\lambda |y,y^{};\overline{q})],\hfill & \mathrm{}=1\hfill \\ g_{\text{MD}}\mu ^2[w_a^{\prime \prime }(\mu |y,y^{};\overline{q})w_a^{\prime \prime }(\lambda |y,y^{};\overline{q})(\lambda ^2\mu ^2)w_b^{\prime \prime }(\lambda |y,y^{};\overline{q})],\hfill & \mathrm{}=2\hfill \end{array}`$ (D63)
$`W_{\text{TN}}^{\prime \prime }(y,y^{};\overline{q})`$ $`=`$ $`\{\begin{array}{cc}g_{\text{TN}}\{[3(yy^{}+\overline{q})^2+\mu ^2]w_a^{\prime \prime }(\mu |y,y^{};\overline{q})\hfill & \\ [(1y^2)/2][(1y_{}^{}{}_{}{}^{2})/2]\},\hfill & \mathrm{}=0\hfill \\ g_{\text{TN}}\{[3(yy^{}+\overline{q})^2+\mu ^2]w_a^{\prime \prime }(\mu |y,y^{};\overline{q})\hfill & \\ [3(yy^{}+\overline{q})^2+\lambda ^2]w_a^{\prime \prime }(\lambda |y,y^{};\overline{q})\},\hfill & \mathrm{}=1\hfill \\ g_{\text{TN}}\{[3(yy^{}+\overline{q})^2+\mu ^2][w_a^{\prime \prime }(\mu |y,y^{};\overline{q})w_a^{\prime \prime }(\lambda |y,y^{};\overline{q})]\hfill & \\ (\lambda ^2\mu ^2)[3(yy^{}+\overline{q})^2+\lambda ^2]w_b^{\prime \prime }(\lambda |y,y^{};\overline{q})\},\hfill & \mathrm{}=2,\hfill \end{array}`$ (D71)
where
$`w_a^{\prime \prime }(\lambda |y,y^{};\overline{q})={\displaystyle \frac{1}{8}}\{`$ (D73)
$`[\lambda ^2+(yy^{}+\overline{q})^2+2y^2y_{}^{}{}_{}{}^{2}]2(2y^2y_{}^{}{}_{}{}^{2})\mathrm{ln}|2[\lambda ^2+(yy^{}+\overline{q})^2]|`$ (D74)
$`+\sqrt{[\lambda ^2+(yy^{}+\overline{q})^2+2y^2y_{}^{}{}_{}{}^{2}]^24(1y^2)(1y_{}^{}{}_{}{}^{2})}`$ (D75)
$`+2(1y^2)`$ (D76)
$`\times \mathrm{ln}\left|\lambda ^2+(yy^{}+\overline{q})^2+y^2y_{}^{}{}_{}{}^{2}+\sqrt{[\lambda ^2+(yy^{}+\overline{q})^2+2y^2y_{}^{}{}_{}{}^{2}]^24(1y^2)(1y_{}^{}{}_{}{}^{2})}\right|`$ (D77)
$`+2(1y_{}^{}{}_{}{}^{2})`$ (D78)
$`\times \mathrm{ln}|\lambda ^2+(yy^{}+\overline{q})^2y^2+y_{}^{}{}_{}{}^{2}+\sqrt{[\lambda ^2+(yy^{}+\overline{q})^2+2y^2y_{}^{}{}_{}{}^{2}]^24(1y^2)(1y_{}^{}{}_{}{}^{2})}|\}`$ (D79)
(D80)
$`w_b^{\prime \prime }(\lambda |y,y^{};\overline{q})={\displaystyle \frac{1}{8}}`$ (D81)
$`\times {\displaystyle \frac{\lambda ^2+(yy^{}+\overline{q})^2+2y^2y_{}^{}{}_{}{}^{2}\sqrt{[\lambda ^2+(yy^{}+\overline{q})^2+2y^2y_{}^{}{}_{}{}^{2}]^24(1y^2)(1y_{}^{}{}_{}{}^{2})}}{\lambda ^2+(yy^{}+\overline{q})^2}}.`$ (D82)
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# Comment on non-Gaussianity in hybrid inflation
## Abstract
In the literature there have been incompatible estimates for the amount of non-Gaussianity in hybrid inflation. In this note we point out the sources for the discrepancies and show that the results for the amount of non-Gaussianity in hybrid inflation obtained by two different methods, namely, perturbing Einstein equation to second order and the separate universe approach, indeed are compatible. This provides confidence in the methods themselves and in the actual computation of non-Gaussianities.
preprint: HIP-2005-26/TH
Introduction. Recently there has been considerable interest in the possible non-Gaussian component of the cosmological perturbations. Non-Gaussian perturbations in hybrid inflation were first estimated using consistent second order perturbation theory<sup>1</sup><sup>1</sup>1This is during inflation. Complete second order perturbation theory connecting inflationary perturbations to the observed CMB anisotropies has so far not been developed. in Enqvist:2004bk , where the formalism of Acquaviva et al. Acquaviva:2002ud was used. The approach is to perturb the metric and matter sides of the Einstein equation, and use the resulting equations to obtain the dynamics of required quantities. These quantities are then used to find out the second order curvature perturbation during inflation.
This method leads to a set of equations and, even after simplifying (although motivated and not too constraining) assumptions, the final expression for the second order curvature is complicated, containing several nonlocal and time-integrated terms. An order of magnitude of the result was estimated already in the original study Enqvist:2004bk . It was later re-estimated in Lyth:2005du , where the result seemed to be much larger.
A completely different method, the $`\delta N`$ formalism or the separate universe approach (see e.g. Wands:2000dp ), for computing the second order curvature perturbation in hybrid inflation was employed by Lyth and Rodríguez in Lyth:2005fi . The result disagrees with the earlier estimates and implies an insignificant contribution from the transverse field, $`\sigma `$, in hybrid inflation.
There is a discrepancy between the results obtained with the two different methods, thus raising doubts on the validity of the methods. However, we argue that in both Enqvist:2004bk and Lyth:2005du the time evolution of certain quantities is not properly taken into account. In this brief comment we present a re-estimate of the original expression for the second order curvature perturbation from the $`\sigma `$ field. We show that when all the time evolutions are correctly taken into account the order of magnitude estimate does indeed agree with the result obtained with the separate universe approach. The second order curvature seems to be proportional to the slow roll parameters. Such a small curvature alone would make the resultant non-Gaussianity unobservable, but according to Boubekeur:2005fj there is a further suppression of the nonlinearity parameter in this particular model due to the uncorrelated Lyth:2005du nature of the $`\sigma `$ non-Gaussianities.
Original computation. The original estimation of the non-Gaussianity was obtained in Enqvist:2004bk by extending the formalism of Acquaviva:2002ud for two scalar fields. The approach is to expand the perturbations of the metric and the matter, which consists of two scalar fields, up to second order. These perturbed quantities are then used to write Einstein equations to second order. The curvature perturbation, defined for one scalar field in the first order as Acquaviva:2002ud $`=\psi +H\frac{\delta \phi }{\dot{\phi }}`$, is written for two scalar fields and expanded to second order; here $`\psi `$ is metric perturbation, $`H=\dot{a}/a`$ is the Hubble parameter, $`\delta \phi `$ is the inflaton perturbation, and $`\dot{\phi }`$ is the time derivative of the background value of the inflaton field. The Einstein equations are then used to obtain the evolution of $``$ to second order, which in turn is then used to estimate the amount of non-Gaussianity.
The analytic calculation of the evolution equations in the case of two scalar fields becomes complicated. To alleviate these difficulties two simplifying assumptions are made in Enqvist:2004bk , namely, it is assumed that the background value $`\sigma _0=0`$ and that the potential does not have any terms linear in $`\sigma `$. The latter assumption means that $`\sigma _0=0`$ indeed is a local minimum. The assumptions are well motivated and not too constraining, and they clearly apply to many other models in addition to hybrid inflation, whose potential is Enqvist:2004bk $`V=V_0\frac{m_0^2}{2}\sigma ^2+\frac{\lambda }{4}\sigma ^4+\frac{m^2}{2}\phi ^2+\frac{g^2}{2}\sigma ^2\phi ^2`$. In fact, at least some number of e-folds after horizon exit of the relevant scales, the form of the potential can be taken to be $`V=V_0+\frac{m_\sigma ^2}{2}\sigma ^2+\frac{m^2}{2}\phi ^2`$, where $`m_\sigma `$ and $`m`$ are the effective masses of $`\sigma `$ and $`\phi `$, respectively.
The two constraints cause the second field, $`\sigma `$, to completely decouple from the first order Einstein equations. Its behaviour in the first order is only governed by the Klein-Gordon equation. At the second order, the contribution of $`\sigma `$ becomes completely additive, i.e. in the evolution equations there are no $`\phi \sigma `$-mixing terms. This enables one to make use of the result of Acquaviva:2002ud for the inflaton contribution in Enqvist:2004bk . Therefore, what is new in Enqvist:2004bk is the contribution coming from the transverse field $`\sigma `$.
The expression obtained in the original paper Enqvist:2004bk for the contribution of the $`\sigma `$ field perturbations to the second order curvature perturbation is rather complicated, containing several nonlocal terms and several time integrated terms. It reads
$`_\sigma ^{\text{(2)}}={\displaystyle \frac{1}{ϵHM_P^2}}\{{\displaystyle }[6H\mathrm{\Delta }^1_i(\delta \dot{\sigma }_1^i\delta \sigma _1)`$
$`+4\mathrm{}^1_i(\delta \dot{\sigma }_1^i\delta \sigma _1)^{\mathbf{}}2(\delta \dot{\sigma }_1)^2+m_\sigma ^2(\delta \sigma _1)^2`$
$`+(ϵ\eta )6H\mathrm{\Delta }^2_i(_k^k\delta \sigma _1^i\delta \sigma _1)^{\mathbf{}}`$
$`+(ϵ\eta )H\mathrm{\Delta }^2_i^i(_k\delta \sigma _1^k\delta \sigma _1)^{\mathbf{}}`$
$`3\mathrm{\Delta }^2_i(_k^k\delta \sigma _1^i\delta \sigma _1)^{\mathbf{}\mathbf{}}`$
$`{\displaystyle \frac{1}{2}}\mathrm{\Delta }^2_i^i(_k\delta \sigma _1^k\delta \sigma _1)^{\mathbf{}\mathbf{}}]dt\mathrm{\Delta }^1_i(\delta \dot{\sigma }_1^i\delta \sigma _1)`$
$`+3\mathrm{\Delta }^2_i(_k^k\delta \sigma _1^i\delta \sigma _1)^{\mathbf{}}+{\displaystyle \frac{1}{2}}\mathrm{\Delta }^2_i^i(_k\delta \sigma _1^k\delta \sigma _1)^{\mathbf{}}`$
$`+3ϵH\mathrm{\Delta }^2_i(_k^k\delta \sigma _1^i\delta \sigma _1)`$
$`+{\displaystyle \frac{ϵH}{2}}\mathrm{\Delta }^2_i^i(_k\delta \sigma _1^k\delta \sigma _1)\},`$ (1)
where $`\mathrm{\Delta }^1`$ is the inverse Laplacian, $`M_P(8\pi G_N)^{1/2}`$ is the reduced Planck mass, $`\delta \sigma _1`$ is the first order perturbation of the transverse scalar field $`\sigma `$, and $`ϵ\frac{M_P^2}{2}\left(\frac{1}{V}\frac{V}{\phi }\right)^2`$ and $`\eta M_P^2\frac{1}{V}\frac{^2V}{\phi ^2}`$ are slow roll parameters; dot denotes derivative with respect to time.
In Enqvist:2004bk the slow roll solution $`\delta \sigma _1e^{m_\sigma ^2t/3H}`$ was used to obtain an estimate for the time derivative $`|\delta \dot{\sigma }_1|\frac{m_\sigma ^2}{H}|\delta \sigma _1|`$. (Double time derivatives were estimated by $`d^2/dt^2(m_\sigma ^2/H)^2`$, but actually one should use equation of motion to get rid of them). Since both fields, $`\phi `$ and $`\sigma `$, are effectively massless, the relation $`|\delta \sigma _1||\delta \phi _1|`$ for the first order perturbations was used to approximate
$$\left|H\frac{\delta \sigma _1}{\dot{\phi }_0}\right|\left|H\frac{\delta \phi _1}{\dot{\phi }_0}\right|\left|^{(1)}\right|.$$
(2)
In the first order Einstein equations $`\sigma `$ is completely decoupled, and the situation is essentially that of a single field inflation. Therefore, $`^{(1)}`$ stays constant outside horizon.
Only order of magnitude estimate was pursued and cancelling spatial derivative operators were neglected, e.g. $`|\mathrm{\Delta }^1_i^{(1)}^i^{(1)}||^{(1)}|^2`$. For the estimation of the time integral the quantities $`H`$, $`ϵ`$, $`\eta `$, $`m_\sigma `$ and $`\delta \sigma _1`$ in the integrand were taken to be constants. The original result in Enqvist:2004bk for the estimate of the second order perturbation due to $`\sigma `$ reads
$$_\sigma ^{(2)}𝒪(ϵ,\eta ,\frac{m_\sigma ^2}{H^2})|^{(1)}|^2.$$
(3)
Since $`\eta _\sigma M_P^2\frac{1}{V}\frac{^2V}{\sigma ^2}\frac{m_\sigma ^2}{H^2}`$, the entire coefficient is of the order slow roll parameters.
Re-estimate by Lyth and Rodríguez. Later Lyth and Rodríguez Lyth:2005du made a re-estimation of Eq. (Comment on non-Gaussianity in hybrid inflation) by inserting initial conditions and writing the equation as a definite integral. They, however, used the same estimates, Eq. (2) (with $`^{(1)}\text{const}`$) and $`|\delta \dot{\sigma }_1|\frac{m_\sigma ^2}{H}|\delta \sigma _1|`$, as the original study Enqvist:2004bk . Similarly, they also assumed $`H`$, $`m_\sigma `$, and $`ϵ`$ to be constants, ending up with<sup>2</sup><sup>2</sup>2They used a different definition for the second order curvature, but that is irrelevant here. Lyth:2005du
$`^{(2)}(t)^{(2)}(t_i)={\displaystyle \frac{1}{ϵHM_P^2}}{\displaystyle _{t_i}^t}[6H\mathrm{\Delta }^1_i(\delta \dot{\sigma }_1^i\delta \sigma _1)`$
$`+m_\sigma ^2(\delta \sigma _1)^22(\delta \dot{\sigma }_1)^2]dt\mathrm{\Delta }N{\displaystyle \frac{m_\sigma ^2}{H^2}}|^{(1)}|^2.`$ (4)
Since $`\mathrm{\Delta }N`$ is the number of e-folds, which can very well be $`60`$, this result implies a much larger $`\sigma `$ contribution to the second order curvature perturbation than in the original study.
Separate universe approach. In addition to the cosmological perturbation theory approach Enqvist:2004bk , there also exists a recent computation Lyth:2005fi of the second order curvature perturbation in hybrid inflation using the separate universe approach.
The general idea of the separate universe approach (see e.g. Wands:2000dp for a concise description) is to consider each point in space as being surrounded by a homogeneous FRW universe. Each point then has its own expansion parameter $`N`$, i.e local number of e-folds, independent of the value of the expansion parameter (or any quantity) in other points. This expansion parameter depends on the values of relevant quantities, such as unperturbed scalar fields, at that point. The complete, inhomogeneous, behaviour of the universe is obtained when all the separately treated points are patched together.
The curvature perturbation, $`\zeta `$, is in Lyth:2005fi defined by
$$g_{ij}=a^2(t)e^{2\zeta (t,𝒙)}\gamma _{ij}(t,𝒙),$$
(5)
where, within inflationary context, $`\gamma _{ij}`$ contains the tensor perturbation which we do not consider here, (see Lyth:2005du for more details). Up to second order in scalar field perturbations ($`\delta \varphi _i\delta \varphi _i(t,𝒙)`$) $`\zeta `$ is obtained from Lyth:2005fi
$$\zeta (t,𝒙)=\underset{i}{}N_{,i}(t)\delta \varphi _i+\frac{1}{2}\underset{ij}{}N_{,ij}(t)\delta \varphi _i\delta \varphi _j,$$
(6)
where $`N_{,i}\frac{N}{\varphi _i}`$ and $`N_{,ij}\frac{^2N}{\varphi _i\varphi _j}`$.
Adapting the notation $`V=V_0(1+\frac{1}{2}\eta \phi ^2+\frac{1}{2}\eta _\sigma \sigma ^2)`$ for the potential, the curvature perturbation reads Lyth:2005fi
$$\zeta =\frac{\delta \phi }{\eta \phi }\frac{\eta }{2}\left(\frac{\delta \phi }{\eta \phi }\right)^2+\frac{\eta _\sigma }{2}e^{2\mathrm{\Delta }N(\eta \eta _\sigma )}\left(\frac{\delta \sigma }{\eta \phi }\right)^2.$$
(7)
The fields $`\phi `$ and $`\sigma `$ are assumed to be massless and their perturbations, $`\delta \phi `$ and $`\delta \sigma `$, are assumed to have the same spectrum $`\left(\frac{H_{}}{2\pi }\right)^2`$ in Lyth:2005fi . Therefore, we can set $`|\zeta _1|\left|\frac{\delta \phi }{\eta \phi }\right|\left|\frac{\delta \sigma }{\eta \phi }\right|`$ and estimate the last term in Eq. (7), i.e. the contribution of $`\sigma `$ to the second order curvature perturbation as
$$\zeta _{2,\sigma }𝒪(\eta _\sigma )e^{2\mathrm{\Delta }N(\eta \eta _\sigma )}|\zeta _1|^2.$$
(8)
Source for the discrepancy. The seeming discrepancy between the two different methods basically comes down to articles Enqvist:2004bk and Lyth:2005du , or more precisely, assuming $`ϵ`$ to be constant, and to the assumption in Eq. (2), i.e.
$$|\delta \sigma _1(t)||\delta \phi _1(t)|,$$
(9)
which is not generally justified. Eq. (9) only holds immediately after horizon exit ($`t=t_i`$), when the amplitude of the perturbation of any effectively massless field $`f`$ is $`|\delta f|H`$.
Using the slow roll equations (outside horizon) we obtain
$`\delta \sigma _1(t)`$ $`=`$ $`\delta \sigma _1(t_i)e^{\eta _\sigma \mathrm{\Delta }N},`$
$`\delta \phi _1(t)`$ $`=`$ $`\delta \phi _1(t_i)e^{\eta _\phi \mathrm{\Delta }N},`$ (10)
$`\phi _0(t)`$ $`=`$ $`\phi _0(t_i)e^{\eta _\phi \mathrm{\Delta }N},`$
where $`\mathrm{\Delta }N=H\mathrm{\Delta }t`$ is the number of e-folds since $`t_i`$; thus, we obtain $`\delta \dot{\sigma }_1=\eta _\sigma H\delta \sigma _1`$, $`\delta \dot{\phi }_1=\eta _\phi H\delta \phi _1`$, and $`\dot{\phi }_0=\eta _\phi H\phi _0`$. Since $`ϵ\dot{\phi }_0^2/H^2M_P^2`$ we can also readily write
$$ϵ(t)=ϵ_ie^{2\eta _\phi \mathrm{\Delta }N},$$
(11)
where we have denoted $`ϵ_iϵ(t_i)`$.
Now, it is immediately clear that $`|^{(1)}|=|H\delta \phi _1/\dot{\phi }_0|`$ stays constant, but one also sees that $`\mathrm{\Delta }N`$ e-folds after horizon exit
$$|H\frac{\delta \sigma _1}{\dot{\phi }_0}||^{(1)}||\frac{\delta \sigma _1}{\delta \phi _1}|e^{\mathrm{\Delta }N(\eta _\phi \eta _\sigma )}|^{(1)}|.$$
(12)
For the estimation of the time integral in Eq. (Comment on non-Gaussianity in hybrid inflation) the important point is that one may not move $`1/ϵ`$ and $`\delta \sigma _1`$ into and out of the time integral.
Re-estimate of non-Gaussianity. Now we present a re-estimate of Eq. (Comment on non-Gaussianity in hybrid inflation) using the time evolutions expressed in Eqs. (Comment on non-Gaussianity in hybrid inflation) and (11). First, we notice that
$`6H\mathrm{\Delta }^1_i(\delta \dot{\sigma }_1^i\delta \sigma _1)+2\mathrm{\Delta }^1_i(\delta \dot{\sigma }_1^i\delta \sigma _1)^{\mathbf{}}`$
$`(\delta \dot{\sigma }_1)^2+m_\sigma ^2(\delta \sigma _1)^2`$
$`=2\mathrm{\Delta }^1_i\left[(3H\delta \dot{\sigma }_1+\delta \ddot{\sigma }_1+m_\sigma ^2\delta \sigma _1)^i\delta \sigma _1\right]`$
$`=0,`$ (13)
since $`3H\delta \dot{\sigma }_1+\delta \ddot{\sigma }_1+m_\sigma ^2\delta \sigma _1=0`$ outside horizon. Thus, Eq. (Comment on non-Gaussianity in hybrid inflation) can be written<sup>3</sup><sup>3</sup>3We use the indefinite integral here, instead of definite one with initial conditions, since the initial second order curvature, or at least initial non-Gaussianity, is supposedly small enough to be safely neglected Lyth:2005du .
$`_\sigma ^{\text{(2)}}`$ $`=`$ $`{\displaystyle \frac{1}{ϵHM_P^2}}\{{\displaystyle ^t}[2\mathrm{}^1_i(\delta \dot{\sigma }_1^i\delta \sigma _1)^{\mathbf{}}(\delta \dot{\sigma }_1)^2`$ (14)
$`+2H(ϵ\eta )\dot{\gamma }_\sigma \ddot{\gamma }_\sigma ]dt\mathrm{\Delta }^1_i(\delta \dot{\sigma }_1^i\delta \sigma _1)`$
$`+\dot{\gamma }_\sigma +ϵH\gamma _\sigma \}`$
$`=`$ $`{\displaystyle \frac{1}{ϵHM_P^2}}\{{\displaystyle ^t}[(\delta \dot{\sigma }_1)^2+2Hϵ\dot{\gamma }_\sigma ]dt`$
$`+\mathrm{\Delta }^1_i(\delta \dot{\sigma }_1^i\delta \sigma _1)+H(ϵ2\eta )\gamma _\sigma \},`$
where we have denoted
$`\gamma _\sigma `$ $``$ $`3\mathrm{\Delta }^2_i(_k^k\delta \sigma _1^i\delta \sigma _1)`$ (15)
$`+`$ $`{\displaystyle \frac{1}{2}}\mathrm{\Delta }^2_i^i(_k\delta \sigma _1^k\delta \sigma _1).`$
The last step is due to time derivatives within the time integral ($`\eta `$ and $`H`$ are assumed to be constants).
Eq. (Comment on non-Gaussianity in hybrid inflation) is used to get rid of the first term in Eq. (Comment on non-Gaussianity in hybrid inflation), $`\frac{1}{ϵHM_P^2}6H\mathrm{\Delta }^1_i(\delta \dot{\sigma }_1^i\delta \sigma _1)dt`$, which would give too large a contribution to the estimate. Note that the order of magnitude estimate does not take into account possible cancellations and, therefore, provides only an upper limit. However, the cancellations can be treated explicitly, as is done here.
For estimation purposes we also adopt the potential used in Lyth:2005fi , $`V=V_0(1+\frac{1}{2}\eta \phi ^2+\frac{1}{2}\eta _\sigma \sigma ^2)`$. The slow roll parameters $`\eta `$, and $`\eta _\sigma `$ are constants, and we also set $`H`$ to be constant. The time evolutions of $`\delta \sigma _1`$, $`\delta \phi _1`$, $`\phi _0`$, and $`ϵ`$ are given by Eqs. (Comment on non-Gaussianity in hybrid inflation) and (11).
Since the order of magnitude estimate anyway gives an upper limit, and since $`ϵϵ_i`$ for any time $`tt_i`$, we replace $`ϵ`$ with $`ϵ_i`$ except in the factor $`\frac{1}{ϵ}`$. We again neglect cancelling orders of spatial derivative operators<sup>4</sup><sup>4</sup>4Because of this, both terms in $`\gamma _\sigma `$ are effectively the same., and put $`\delta \dot{\sigma }_1=\eta _\sigma H\delta \sigma _1`$. The estimate for the second order curvature perturbation, Eq. (14), thus becomes
$`_\sigma ^{\text{(2)}}`$ $``$ $`{\displaystyle \frac{1}{ϵHM_P^2}}\{{\displaystyle ^t}[\eta _\sigma ^2H^2|\delta \sigma _1|^2+ϵ_i\eta _\sigma H^2|\delta \sigma _1|^2]dt`$ (16)
$`+\eta _\sigma H|\delta \sigma _1|^2+ϵ_iH|\delta \sigma _1|^2+\eta H|\delta \sigma _1|^2\}`$
$``$ $`{\displaystyle \frac{1}{ϵHM_P^2}}\{𝒪(ϵ_i,\eta _\sigma ){\displaystyle ^t}\eta _\sigma H^2|\delta \sigma _1|^2dt`$
$`+𝒪(ϵ_i,\eta ,\eta _\sigma )H|\delta \sigma _1|^2\}.`$
We have now two terms to evaluate, namely
$`{\displaystyle \frac{1}{ϵHM_P^2}}H|\delta \sigma _1|^2={\displaystyle \frac{|\delta \sigma _1(t_i)|^2}{ϵ_iM_P^2}}e^{2\mathrm{\Delta }N(\eta \eta _\sigma )}`$
$`=\left|H{\displaystyle \frac{\delta \sigma _1(t_i)}{\dot{\phi }_0(t_i)}}\right|^2e^{2\mathrm{\Delta }N(\eta \eta _\sigma )}`$
$`=e^{2\mathrm{\Delta }N(\eta \eta _\sigma )}|^{(1)}|^2,`$ (17)
and
$`{\displaystyle \frac{1}{ϵHM_P^2}}{\displaystyle ^t}H^2\eta _\sigma |\delta \sigma _1|^2𝑑t`$
$`={\displaystyle \frac{|\delta \sigma _1(t_i)|^2}{ϵM_P^2}}\eta _\sigma {\displaystyle ^{\mathrm{\Delta }N}}e^{2\eta _\sigma N}𝑑N`$
$`{\displaystyle \frac{|\delta \sigma _1(t_i)|^2}{ϵM_P^2}}e^{2\eta _\sigma \mathrm{\Delta }N}=e^{2\mathrm{\Delta }N(\eta \eta _\sigma )}|^{(1)}|^2.`$ (18)
Therefore, our final estimate reads
$$_\sigma ^{\text{(2)}}𝒪(ϵ_i,\eta ,\eta _\sigma )e^{2\mathrm{\Delta }N(\eta \eta _\sigma )}|^{(1)}|^2.$$
(19)
Comparing to the estimate in Eq. (8) one sees that the two results are of the same form and of the same order. Instead of $`𝒪(ϵ_i,\eta ,\eta _\sigma )`$ Eq. (8) only has a factor $`𝒪(\eta _\sigma )`$. Here we have, however, provided only an order of magnitude estimate and there still is a possibility for cancellations which may have been overlooked. It may be worth pointing out that if the nonlocal terms are discarded in Eq. (14) and only the first term is estimated, we obtain an $`𝒪(\eta _\sigma )`$ coefficient only.
Discussion. The maximum number of e-folds for the observable scales is $`60`$ and observational limits for the spectral index of the curvature perturbation require $`|\eta |0.01`$. It is therefore unlikely that the exponential factor $`e^{2\mathrm{\Delta }N(\eta \eta _\sigma )}`$ would provide any significant enhancement, and the overall factor of the order slow roll parameters gives the magnitude of the result.
The quantity measured in CMB experiments is the nonlinearity parameter $`f_{\text{NL}}`$, which is related to $`^{(2)}`$, but according to Boubekeur:2005fj it is highly suppressed in this scenario. The projected sensitivity (using first order perturbation theory) for an ideal experiment is no better than $`f_{\text{NL}}1`$ even including polarisation Babich:2004yc . It was later realized by Vernizzi Vernizzi:2004nc that the the second order curvature perturbation defined in Acquaviva:2002ud has an artificial time evolution $`\dot{}^{(2)}(2\dot{ϵ}\dot{\eta })_{}^{(1)}{}_{}{}^{2}`$. However, even taking this small effect into account, it seems safe to say that the non-Gaussianity produced in the hybrid scenario by the $`\sigma `$ field seems to be too small to be observed unless some key aspects of the scenario are changed.
There are still conceptual and practical problems with computing non-Gaussianities. The second order theory is not yet well established and the connection between theoretical calculations and CMB observations is far from complete. One problem is that the usual scalar-vector-tensor decomposition of the perturbations is inherently non-local Ellis:1990gia . These non-localities do not appear in the first order, but in the second order equations there are terms like $`\mathrm{\Delta }^1(_ig^ig)`$ and $`\mathrm{\Delta }^1(g\mathrm{\Delta }g)`$, where $`g`$ represents a generic perturbation. The physical interpretation of these terms is not clear.
Lyth and Rodríguez apply the separate universe approach in Lyth:2005fi to compute the second order perturbations, or non-Gaussianity, in various scenarios including hybrid inflation. The formalism they use does not produce non-local terms involving the inverse Laplacian $`\mathrm{\Delta }^1`$, and they state that Lyth:2005fi “…such terms must cancel if correctly evaluated.” Needless to say, this is quite a strong claim. It is actually not certain whether the separate universe approach is completely correct when expanded to second order in perturbations. Indeed, in Wands:2000dp it was stated that any nonlinear interaction introduces mode-mode couplings which undermine the separate universe picture. As an example, recent studies of non-Gaussianities in preheating Enqvist:2004ey demonstrate that these mode-mode couplings seem to be important.
Despite all the problems and ambiguities in the different second order formalisms it is comforting that the two different approaches discussed here can now be seen to produce the same result for the second order curvature perturbation and, therefore, for the amount of non-Gaussianity in hybrid inflation.
Acknowledgements. The author would like to thank Kari Enqvist and Asko Jokinen for useful discussions and commenting on the manuscript and Filippo Vernizzi for enlightening discussions. This work is supported by the Magnus Ehrnrooth Foundation.
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# Translationally invariant discrete kinks from one-dimensional maps
## Abstract
For most discretisations of the $`\varphi ^4`$ theory, the stationary kink can only be centered either on a lattice site or midway between two adjacent sites. We search for exceptional discretisations which allow stationary kinks to be centered anywhere between the sites. We show that this translational invariance of the kink implies the existence of an underlying one-dimensional map $`\varphi _{n+1}=F(\varphi _n)`$. A simple algorithm based on this observation generates three different families of exceptional discretisations.
Since the early 1960s, the $`\varphi ^4`$-equation,
$$\varphi _{tt}=\varphi _{xx}+\frac{1}{2}\varphi (1\varphi ^2),$$
(1)
has been one of the workhorses of statistical mechanics Bishop\_Schneider and quantum field theory Rajaraman . Its kink solution,
$$\varphi (x,t)=\mathrm{tanh}\frac{xctx^{(0)}}{2\sqrt{1c^2}},$$
(2)
together with the sine-Gordon kink, are the simplest examples of topological solitons. More recently, interest has shifted towards the discrete $`\varphi ^4`$-theories Sanchez ; phi4kinks ,
$$\ddot{\varphi }_n=\frac{\varphi _{n+1}2\varphi _n+\varphi _{n1}}{h^2}+f(\varphi _{n1},\varphi _n,\varphi _{n+1}),$$
(3)
and their solutions. Here $`h`$ is the lattice spacing: $`\varphi _n(t)=\varphi (x_n,t)`$, with $`x_n=hn`$, and the function $`f`$ is chosen to reproduce the nonlinearity (1) in the continuum limit:
$$f(\varphi ,\varphi ,\varphi )=\frac{1}{2}\varphi (1\varphi ^2).$$
(4)
For a variety of discrete nonlinearities $`f`$, Eq.(3) admits stationary kink solutions Slovenians ; Yip . The discrete $`\varphi ^4`$ kinks have been used to describe incommensurate systems and narrow domain walls in ferroelectrics and ferromagnets, topological excitations in biological macromolecules and hydrogen-bonded chains, and polymerization mismatches in polymers Sanchez ; Applications . Physically, one of the most significant properties of domain walls and topological defects is their mobility Yip ; however whether the discrete equations (3) can admit travelling kinks remains an open question phi4kinks ; moving\_kinks ; Speight1 ; Speight2 ; Speight3 ; Panos . The continuous $`\varphi ^4`$-equation (1) is Lorentz-invariant, and so the existence of the travelling kink (2) is an immediate consequence of the existence of the stationary soliton, Eq.(2) with $`c=0`$. The discretisation breaks the Lorentz invariance and the existence of travelling discrete kinks becomes a nontrivial matter.
In fact, the discretisation even breaks the translation invariance of Eq.(1). As a result, the stationary kink can be centered only at a countable number of points — usually on a site and midway between two adjacent sites Slovenians ; Yip . This breaking of the translation invariance is connected with the presence of the Peierls-Nabarro barrier, an additional periodic potential induced by discreteness.
Miraculously, there are several exceptional discretisations which, while breaking the translation invariance of the equation, allow the existence of translationally invariant kinks; that is, kinks centred at an arbitrary point between the sites. One such discretisation was discovered by Speight and Ward using a Bogomolny-type energy-minimality argument Speight1 ; Speight2 :
$`f={\displaystyle \frac{2\varphi _n+\varphi _{n+1}}{12}}\left(1{\displaystyle \frac{\varphi _n^2+\varphi _n\varphi _{n+1}+\varphi _{n+1}^2}{3}}\right)`$
$`+{\displaystyle \frac{2\varphi _n+\varphi _{n1}}{12}}\left(1{\displaystyle \frac{\varphi _n^2+\varphi _n\varphi _{n1}+\varphi _{n1}^2}{3}}\right).`$ (5)
Another one derives from the Ablowitz-Ladik integrable discretisation of the nonlinear Schrödinger equation; it was reobtained by Bender and Tovbis BT97 from the requirement of suppression of the kinks’ resonant radiation:
$$f=\frac{1}{4}\left(\varphi _{n+1}+\varphi _{n1}\right)\left(1\varphi _n^2\right).$$
(6)
Finally, the nonlinearity
$`f={\displaystyle \frac{\varphi _{n+1}+\varphi _{n1}}{4}}{\displaystyle \frac{(\varphi _{n+1}^2+\varphi _{n1}^2)(\varphi _{n+1}+\varphi _{n1})}{8}}`$ (7)
was identified by Kevrekidis Panos , who demonstrated the existence of a two-point invariant associated with the stationary equation
$$\frac{\varphi _{n+1}2\varphi _n+\varphi _{n1}}{h^2}+f(\varphi _{n1},\varphi _n,\varphi _{n+1})=0,$$
(8)
with $`f`$ as in (6) and (7).
Although the translation invariance of a stationary kink does not automatically guarantee the existence of a travelling soliton, it is natural to expect it to be a prerequisite for kink mobility. For example, in the variational description of the slowly moving kink, the solution is sought as a stationary kink with a free continuous parameter defining its position on the line Speight1 . Also, the Stokes constants measuring the intensity of resonant radiation from the translationally invariant kinks were found to be at least an order of magnitude smaller than the corresponding constants in models with noninvariant kinks ODI . With an eye to a future attack on travelling kinks, it would be useful to identify all discretisations of the $`\varphi ^4`$ theory supporting translationally invariant stationary kinks. The purpose of this note is to provide a general recipe for the generation of such exceptional discretisations $`f(\varphi _{n1},\varphi _n,\varphi _{n+1})`$.
We start with a simple observation which, however, holds the key to our construction. Assume we have a nonlinearity $`f`$ which supports a stationary discrete kink, i.e. a monotonically growing sequence $`\varphi _n`$: $`1<\varphi _n<\varphi _{n+1}<1`$ for $`\mathrm{}<n<\mathrm{}`$. Furthermore, assume there exists a continuous monotonically growing function $`g(x)`$, defined for all real $`x`$, such that $`\varphi _n=g(n)`$. (The function $`g`$ can also depend on $`h`$ parametrically but we omit this dependence for simplicity of notation.) This function generates a family of kinks centered at an arbitrary point $`x^{(0)}`$ on the line: $`\varphi _n=g(nx^{(0)})`$. It is important to emphasise that such a continuous function can exist only for exceptional discretisations. For generic discretisations, the function $`g(x)`$ can be defined only on integers.
The existence of the function $`g(x)`$ defined on the entire real line — or, equivalently, the translation invariance of the kink — implies that the stationary equation (8) derives from a two-point map. Indeed, since $`g(x)`$ is monotonic, we can write $`n=g^1(\varphi _n)`$. Now since $`\varphi _{n+1}=g(n+1)`$, we have $`\varphi _{n+1}=g(g^1(\varphi _n)+1)F(\varphi _n)`$, which is a well-defined one-dimensional map.
This observation suggests the following strategy for the construction of exceptional discretisations. Assume we have a 1D map which we will write in the form
$$\varphi _{n+1}\varphi _n=hH(\varphi _{n+1},\varphi _n).$$
(9)
Let $`H`$ satisfy the following continuity condition:
$$H(\varphi ,\varphi )=\frac{1}{2}(1\varphi ^2).$$
(10)
This condition is necessary to make sure that the map (9) becomes
$$\varphi _x=\frac{1}{2}(1\varphi ^2)$$
(11)
in the continuum limit. The stationary ($`c=0`$) kink solution (2) of Eq.(1) is, simultaneously, a solution of the first-order equation (11). Imposing (10) we ensure that the discrete kink of (9) will have the correct continuum limit. Next, Eq.(10) implies that the map (9) has just one pair of fixed points, $`\varphi _{}=1`$. For small $`h`$, $`\varphi _{n+1}`$ remains close to $`\varphi _n`$ and hence, $`H(\varphi _{n+1},\varphi _n)`$ remains close to (10) which is positive for $`|\varphi |<1`$. Consequently, no matter what $`|\varphi _0|<1`$ we start with, the sequence $`\varphi _n`$ is monotonically growing — at least until $`|\varphi _n|`$ is not very close to 1. To ensure that it remains such near the fixed points, we assume that $`\varphi _{}=1`$ is a source and $`\varphi _{}=1`$ a sink. (That is, small perturbations $`\delta \varphi _n=\varphi _n\varphi _{}`$ satisfy $`\delta \varphi _{n+1}=\lambda \delta \varphi _n`$ with $`\lambda >1`$ near $`\varphi _{}=1`$ and $`0<\lambda <1`$ near $`\varphi _{}=1`$.) Then, for any $`h`$ smaller than some $`\overline{h}`$ and any $`\varphi _0`$ between $`1`$ and $`1`$, there is a number $`N`$ such that $`|\varphi _N\varphi _{}|`$ is so small that all $`\varphi _n`$ with $`n>N`$ are entrapped by the “linear neighbourhood” of $`\varphi _{}=1`$ and those with $`n<N`$ are all in a neighbourhood of $`\varphi _{}=1`$. This means that each $`\varphi _0`$ with $`|\varphi _0|<1`$ defines a monotonic kink solution and so for any sufficiently small $`h`$ we have a one-parameter family of stationary kinks. Speight Speight3 gives a less intuitively appealing but more rigorous proof of this fact; he also shows that our assumption on the character of the fixed points can be relaxed.
Next, squaring both sides of (9) and subtracting the square of its back-iterated copy,
$$\varphi _n\varphi _{n1}=hH(\varphi _n,\varphi _{n1}),$$
(12)
produces an exceptional stationary Klein-Gordon equation
$$\frac{\varphi _{n+1}2\varphi _n+\varphi _{n1}}{h^2}=\frac{H^2(\varphi _{n+1},\varphi _n)H^2(\varphi _n,\varphi _{n1})}{\varphi _{n+1}\varphi _{n1}}.$$
(13)
If $`H`$ is symmetric: $`H(\varphi _n,\varphi _{n1})=H(\varphi _{n1},\varphi _n)`$, the numerator vanishes exactly where the denominator equals zero, so the discretisation (13) is nonsingular.
If we want to have polynomial discretisations of the $`\varphi ^4`$ theory, the function $`H^2`$ has to be a quartic polynomial. This leads to two possibilities, one where $`H`$ is the square root of a polynomial, and the other where $`H`$ is a polynomial itself. These can be written jointly as
$$\left(\varphi _{n+1}\varphi _n\right)^m=h^mP_{2m}(\varphi _{n+1},\varphi _n),$$
(14)
where $`m=1`$ or $`2`$, and $`P_{2m}(u,v)`$ is a polynomial of degree $`2m`$ that satisfies the symmetry and continuity conditions
$`P_{2m}(u,v)`$ $`=`$ $`P_{2m}(v,u),`$ (15)
$`P_{2m}(\varphi ,\varphi )`$ $`=`$ $`2^m(1\varphi ^2)^m.`$ (16)
The condition (16) is a consequence of Eq.(10).
Before we proceed to the classification of the resulting models, it is pertinent to note that the linear part of the function $`f`$ in (3) can always be fixed to $`\frac{1}{2}\varphi _n`$ without loss of generality. Indeed, the most general function satisfying (4) is $`f=a\varphi _n+\frac{1}{2}\left(\frac{1}{2}a\right)\left(\varphi _{n+1}+\varphi _{n1}\right)+\text{cubic terms}`$. Since $`h^2`$ in (8) is a free parameter, we can always make a replacement $`h\stackrel{~}{h}`$ such that $`a2/h^2=\stackrel{~}{a}2/\stackrel{~}{h}^2`$. In particular, we can set $`\stackrel{~}{a}=\frac{1}{2}`$ which gives
$$f(\varphi _{n1},\varphi _n,\varphi _{n+1})=\frac{1}{2}\varphi _nQ(\varphi _{n1},\varphi _n,\varphi _{n+1}),$$
(17)
where $`Q`$ is a homogeneous polynomial of degree 3.
Let, now, $`m=2`$ in Eq.(14). Provided $`P_4`$ satisfies conditions (15) and (16), the numerator $`P_4(\varphi _{n+1},\varphi _n)P_4(\varphi _n,\varphi _{n1})`$ of the fraction in the right-hand side of Eq.(13) divides $`(\varphi _{n+1}\varphi _{n1})`$ and so Eq.(13) will be of the form (8) with some cubic function $`f`$. The most general choice for such a polynomial is
$$\begin{array}{c}P_4(u,v)=\frac{1}{4}\mu (uv)^2\frac{1}{2}uv\hfill \\ \hfill +\frac{1}{20}\left[\alpha (u^4+v^4)+\beta uv(u^2+v^2)+\gamma u^2v^2\right],\end{array}$$
(18)
where $`\alpha `$, $`\beta `$, $`\gamma `$ satisfy $`2\alpha +2\beta +\gamma =5`$ and $`\mu `$ is arbitrary. Picking the positive value of $`\sqrt{P_4}`$ and assuming that $`h`$ is sufficiently small, one can check that the fixed points $`\varphi _{}=1`$ of the map (14) are a source and a sink, for any $`\mu ,\alpha `$ and $`\beta `$. Consequently, the resulting cubic polynomial,
$$\begin{array}{c}Q=\frac{1}{20}[\alpha (\varphi _{n+1}+\varphi _{n1})(\varphi _{n+1}^2+\varphi _{n1}^2)+\gamma \varphi _n^2(\varphi _{n+1}\hfill \\ \hfill +\varphi _{n1})+\beta \varphi _n(\varphi _{n+1}^2+\varphi _n^2+\varphi _{n1}^2+\varphi _{n+1}\varphi _{n1})]\end{array}$$
(19)
with $`\gamma =52(\alpha +\beta )`$ defines a two-parameter family of models with translationally invariant kink solutions.
The discretisation (19) includes, as particular cases, the Bender-Tovbis function (6) (which results from setting $`\alpha =\beta =0`$) and the Kevrekidis nonlinearity (7) (for which $`\alpha =\frac{5}{2}`$, $`\beta =0`$). Another simple function arises by letting $`\alpha =\gamma =0`$; this is a new model:
$$Q=\frac{1}{8}\varphi _n(\varphi _{n+1}^2+\varphi _n^2+\varphi _{n1}^2+\varphi _{n+1}\varphi _{n1}).$$
Now let $`m=1`$. The most general quadratic $`P_2`$ satisfying (15)–(16) is
$$P_2(u,v)=\frac{1}{2}\alpha (uv)^2\frac{1}{2}uv,$$
(20)
with an arbitrary $`\alpha `$. Note that for $`h<2`$ and any $`\alpha `$, $`\varphi _{}=1`$ are a source and a sink. Hence all resulting models will exhibit continuous families of kinks. Substituting Eq.(20) for $`H`$ in (13), we obtain just a particular case of the nonlinearity (19), corresponding to the choice of the quartic (18) in the form of a complete square: $`P_4=P_2^2`$. To obtain new models, we need to note a symmetry $`I_3(\varphi _{n1},\varphi _n,\varphi _{n+1})=0`$ which follows from Eq.(14) with $`m=1`$. Here
$$I_3P_2(\varphi _{n+1},\varphi _n)(\varphi _n\varphi _{n1})+P_2(\varphi _n,\varphi _{n1})(\varphi _n\varphi _{n+1}).$$
Equation (13) remains valid if $`\beta I_3`$ is added to its right-hand side, with an arbitrary coefficient $`\beta `$. The resulting function $`Q`$ has the form
$$\begin{array}{c}Q=\alpha ^2(\varphi _{n+1}^3+\varphi _{n1}^3)+2\gamma \left(\alpha \beta \right)\varphi _{n+1}\varphi _n\varphi _{n1}\hfill \\ \hfill +\alpha \left(\alpha \beta \right)\varphi _{n+1}\varphi _{n1}(\varphi _{n+1}+\varphi _{n1})\\ \hfill +\left[2\alpha ^2+\gamma ^2+\beta (\gamma \alpha )\right]\varphi _n^2(\varphi _{n+1}+\varphi _{n1})\\ \hfill +\alpha \left(2\gamma +\beta \right)\varphi _n(\varphi _{n+1}^2+\varphi _{n1}^2)+2\alpha \left(\gamma +\beta \right)\varphi _n^3,\end{array}$$
(21)
where $`\gamma =\frac{1}{2}2\alpha `$. Eq.(21) defines a two-parameter family of discretisations supporting translationally-invariant kinks. These models cannot be obtained within the $`m=2`$ analysis above — unless $`\beta =0`$, of course.
Letting $`\alpha =\beta =\frac{1}{6}`$, we recover the model of Speight and Ward, Eq.(5). Another particularly simple, new, model is obtained by taking $`\alpha =0`$ and $`\beta =\frac{1}{2}`$:
$$Q=\frac{1}{2}\varphi _{n1}\varphi _n\varphi _{n+1}.$$
It is instructive to compare discretisations furnished by our method with those arising from the requirement of the absence of resonant radiation from the kink BT97 . The advance-delay equation associated with Eq.(8),
$$\begin{array}{c}\varphi (x+h)2\varphi (x)+\varphi (xh)\hfill \\ \hfill +h^2f(\varphi (xh),\varphi (x),\varphi (x+h))=0,\end{array}$$
(22)
can be solved to all orders as a perturbation expansion in $`h`$ ODI ; the resultant solution depends continuously on the position parameter $`x^{(0)}`$. It is, therefore, only the terms which lie beyond all orders in $`h`$ which present an obstacle to the existence of a translationally invariant kink. These terms vanish if Stokes constants at all orders vanish. The leading Stokes constant vanishes if there exists a convergent solution in powers of $`z^1`$ to the equation
$$\begin{array}{c}\phi (z+1)2\phi (z)+\phi (z1)\hfill \\ \hfill Q(\phi (z1),\phi (z),\phi (z+1))|_{h=0}=0\end{array}$$
(23)
(see e.g. Tovbis1 ). This equation comes from a rescaling of Eq.(22) near the singularities of its leading-order solution (2) at $`x_n=\pi i(1+2n)`$, $`n`$, in the limit $`h0`$. The convergence of a power-series solution to Eq.(23) is necessary for the absence of oscillatory radiation tails in its “parent” equation (22).
In general, a numerical procedure is required to determine whether the series converges for a given $`Q`$, but we can easily generate a class of models for which it truncates after the first term. This was the method employed in Ref.BT97 in deriving Eq.(6). It is a matter of direct substitution to check that the most general cubic polynomial for which $`\phi =2/z`$ is a solution of Eq.(23), is
$$\begin{array}{c}Q=\sigma \varphi _n(\varphi _{n+1}+\varphi _{n1})^22\sigma \varphi _{n+1}\varphi _{n1}(\varphi _{n+1}+\varphi _{n1})\hfill \\ \hfill +\left(\frac{1}{4}\frac{\beta }{2}\right)\varphi _n^2(\varphi _{n+1}+\varphi _{n1})+\beta \varphi _{n1}\varphi _n\varphi _{n+1},\end{array}$$
(24)
with $`\sigma ,\beta `$ arbitrary constants.
The fact that the leading-order Stokes constant is zero is necessary but not sufficient for equation (22) to have continuous families of kinks for finite $`h`$. We tested, numerically, a particular representative from the class (24):
$$Q=\varphi _{n+1}\varphi _{n1}\frac{\varphi _{n+1}+\varphi _{n1}}{2}\varphi _n\frac{\varphi _{n+1}^2+\varphi _{n1}^2}{4}.$$
(25)
This model is obtained by letting $`\sigma =\frac{1}{4}`$ and $`\beta =\frac{1}{2}`$. To check whether translationally-invariant kinks exist or not, we have computed a stationary on-site kink for the model (25) and calculated eigenvalues of the associated linearised operator for an equidistant sequence of $`h`$-values, ranging from $`h=0`$ to $`h=1.179`$ with increment $`0.001`$. For $`h`$ smaller than $`0.556`$, the smallest-modulus eigenvalue was found to be smaller than $`10^{12}`$ which is our numerical error of computation. However as $`h`$ increases from $`0.556`$, the smallest eigenvalue grows to $`\lambda =9\times 10^7`$ at $`h=0.955`$, then decreases, crosses through zero at $`h=0.993`$, after which grows in modulus to $`\lambda =6\times 10^4`$ at $`h=1.179`$ (Fig.1). Thus, the zero eigenvalue, indicating the existence of a continuous family of solutions, is not present in the spectrum for $`h>0.556`$. For $`h<0.556`$, the smallest $`\lambda `$ is apparently also nonzero (though very small). The only exception is the value $`h=0.993`$ for which the zero mode does exist. This means that the model (25) supports a continuous family of kinks for just one, isolated, value of $`h`$.
However, there exists a family of exceptional discretisations which reduces to (24) in the limit $`h0`$. Indeed, when $`\alpha =0`$ in Eq.(20), the map (14) has one more symmetry: $`I_2(\varphi _{n1},\varphi _n,\varphi _{n+1})=0`$, where
$$I_2\left(1+\frac{h^2}{4}\right)\varphi _n(\varphi _{n+1}+\varphi _{n1})2\varphi _{n1}\varphi _{n+1}\frac{h^2}{2}.$$
We can add $`\sigma (\varphi _{n+1}+\varphi _{n1})I_2`$ to the right-hand side of (13), along with $`\beta I_3`$ (with $`\sigma ,\beta `$ arbitrary constants.) This gives rise to the following family of discretisations:
$$\begin{array}{c}Q=\sigma \varphi _n(\varphi _{n+1}+\varphi _{n1})^22\sigma \varphi _{n1}\varphi _{n+1}(\varphi _{n+1}+\varphi _{n1})\hfill \\ \hfill +\left(\frac{1}{4}\frac{\beta }{2}\right)\varphi _n^2(\varphi _{n+1}+\varphi _{n1})+\beta \varphi _{n1}\varphi _n\varphi _{n+1}\\ \hfill +\frac{1}{4}\sigma h^2\varphi _n(\varphi _{n+1}+\varphi _{n1})^2.\end{array}$$
(26)
Except for the last, $`𝒪(h^2)`$, term, this coincides with Eq.(24).
For the map (14)+(20) with $`\alpha =0`$, the kink can be found explicitly. This implies that the discretisations (26) also share an explicit kink solution (for all $`\beta `$ and $`\sigma `$): $`\varphi _n=\mathrm{tanh}(anx^{(0)})`$, with $`\mathrm{tanh}a=h/2`$.
Our final remark is on the conserved quantities of Eq.(3). The translation invariance of the stationary kink does not imply the invariance of equation (3) and hence the conservation of momentum. The discretisation (13) \[and hence (19)\] conserves momentum Panos whereas the nonlinearities (21) and (26) (with $`\beta ,\sigma 0`$) — do not. Moreover, that the discretisation $`f`$ is exceptional does not guarantee that Eq.(3) has any integral of motion whatsoever. In particular, out of the three families (19), (21), and (26), only one model conserves energy, namely Speight and Ward’s, Eq.(5).
In conclusion, we have identified three families of discretisations of the $`\varphi ^4`$ equation which support translationally invariant stationary kinks: Eqs.(19), (21) and (26). In each case we have exhibited, explicitly, the underlying 1D map.
###### Acknowledgements.
I.B. is a Harry Oppenheimer Fellow. O.O. was supported by a Nelli Brown Spilhaus scholarship.
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# Kaluza-Klein 5D Ideas Made Fully Geometric
## 1 Physics: 5D Geometry is Useful
After the 1916 success of A. Einstein, who explained gravitation by combining space and time into a 4D space, there have been many efforts to explain other physical fields by adding other physical dimensions.
The first successful attempt was made by Th. Kaluza and O. Klein in 1921. They showed that if we formally consider the equations of general relativity theory in the 5D space, the equations for the normal $`4\times 4`$ components $`g_{ij}`$ of the metric tensor still describe gravitation, while the new components $`g^{5i}`$ of the metric tensor satisfy Maxwell’s equations (under the assumption that $`g_{55}=\mathrm{const}`$). Thus, if we go to 5D space, we get a geometric interpretation of electrodynamics.
The only problem with this interpretation is that it is formal: change in first 4 dimensions makes perfect physical sense, while there seemed to be no physical effects corresponding to change in 5th dimension. To solve this problem, A. Einstein and P. Bergmann proposed, in 1938 \[Einstein and Bergmann 1938\], that the 5th dimension forms a tiny circle, so that only micro-particles “see” it, while for us, the world is 4D.
This is a standard view now in particle physics; see, e.g., \[Green et al. 1988, Polchinski 1998\]: space is 10- or 11-dimensional, all dimensions except the first four are tiny.
## 2 Formulas from Physical 5D Theories that Need to Be Explained in Purely Geometric Terms
In addition to a nice geometric model, the traditional description of Kaluza-Klein theory requires several additional physical formulas, formulas that look very artificial because they do not have a direct geometric explanation.
In this paper, we will show that, if we take the Einstein-Bergmann model seriously, then these formulas can be derived – and thus, they are not additional and ad hoc.
What are these formulas that do not directly follow from the geometric model?
First, the assumption $`g_{55}=\mathrm{const}`$ is artificial.
Second, since only four coordinates have a physical sense, the distance $`\mathrm{\Delta }s^2=\underset{i=1}{\overset{5}{}}\underset{j=1}{\overset{5}{}}g_{ij}\mathrm{\Delta }x_i\mathrm{\Delta }x_j`$ between the points $`x`$ and $`x+\mathrm{\Delta }x`$ should only depend on the first 4 coordinates – while in general, for a 5D metric, the terms $`g_{55}(\mathrm{\Delta }x^5)^2`$ and $`g_{5i}\mathrm{\Delta }x^5\mathrm{\Delta }x^i`$ create a difficult-to-explain dependence on $`\mathrm{\Delta }x^5`$.
Third, we would like to explain the fact that the observed values of physical fields do not depend on the fifth coordinate $`x^5`$, e.g., that $`g_{ij}/x^5=0`$ (this condition is called cylindricity).
Several other formulas came from the attempts to give the fifth dimension a physical interpretation. Namely, in the 1940s, Yu. Rumer showed (see, e.g., \[Rumer 1956\]) that if we interpret $`x^5`$ as action $`S=L𝑑x𝑑t`$ (i.e., the quantity whose extrema define the field’s dynamics), then the fact that $`x^5`$ is defined on a circle is consistent with the fact that in quantum physics (e.g., in its Feynman integral formulation), action is used only as part of the expression $`\mathrm{exp}(iS/h)`$, whose value is not changed if we add a constant $`2\pi h`$ to $`S`$. (For a H atom, this idea leads to the original Bohr’s quantization rules.)
Action is defined modulo arbitrary transformation $`SS+f(x^i)`$; thus, the corresponding transformation $`x^5x^5+f(x^i)`$ should be geometrically meaningful. Similar transformations stem from the electrodynamic interpretation of $`g_{5i}`$ as $`A_i`$: gauge transformations $`A_iA_if/x_i`$.
## 3 Natural Idea and Its Problems
The main difference between a standard 4D space and Einstein-Bergmann’s 5D model is that we have a cylinder $`K=R^4\times S^1`$ ($`K`$ for Kaluza) instead of a linear space. It is, therefore, desirable to modify standard geometry by substituting $`K`$ instead of $`R^4`$ into all definitions.
The problem with this idea is that the corresponding formalisms of differential geometry use the underlying linear space structure, i.e., addition and multiplication by a scalar. We still have addition in $`K`$, but multiplication is not uniquely defined for angle-valued variables: we can always interpret an angle as a real number modulo the circumference, but then, e.g., $`02\pi `$ while $`0.60\sim ̸0.62\pi `$.
## 4 What We Suggest
We do need a real-number representation of an angle variable. A more natural representation of this variable is not as a single real number, but as a set $`\{\alpha +n2\pi \}`$ of all possible real numbers that correspond to the given angle.
Similarly to interval and fuzzy arithmetic, we can naturally define element-wise arithmetic operations on such sets, e.g., $`A+B=\{a+b|aA,bB\}`$. We can then define tensors as linear mappings that preserve the structure of such sets, and we can define a differentiable tensor field as a field for which the set of all possible values of the corresponding partial derivatives is also consistent with the basic structure.
Comment. These results were first announced in \[Kreinovich and Nguyen 2005, Kreinovich and Starks 1997, Starks and Kreinovich 1998\].
## 5 Resulting Formalism: Idea
In mathematical terms, the resulting formalism is equivalent to the following: We start with the space $`K`$ which is not a vector space (only an Abelian group). We reformulate standard definitions of vector and tensor algebra and tensor analysis and apply them to $`K`$: $`K`$-vectors are defined as elements of $`K`$; $`K`$-covectors as elements of the dual group, etc. All physically motivated conditions turn out to be natural consequences of this formalism.
## 6 $`K`$-Vectors
In the traditional 4-D space-time $`R^4`$, we can define a vector as simply an element of $`R^4`$. In our case, instead of 4-D space-time $`R^4`$, we have a 5-D space-time $`K\stackrel{\mathrm{def}}{=}R^4\times S^1`$, in which $`S^1`$ is a circle of a small circumference $`h>0`$ – i.e., equivalently, a real line in which two numbers differing by a multiple of $`h`$ describe the same point: $`(x^1,\mathrm{},x^4,x^5)(x^1,\mathrm{},x^4,x^5+kh)`$. Thus, it is natural to define $`K`$-vectors as simply elements of $`K`$:
Definition 1. A $`K`$-vector is an element of $`K=R^4\times S^1`$.
On the set of all vectors in $`R^4`$, there are two natural operations: (commutative) addition $`a+b`$ and multiplication by a real number $`\lambda `$: $`a\lambda a`$. Thus, this set is a linear space.
In contrast, on the the set $`K`$ of all $`K`$-vectors we only have addition, so the set of all $`K`$-vectors is not a linear space, it is only an Abelian group.
## 7 $`K`$-Covectors
In physics, an important algebraic object is a covector: vectors describe the location $`x`$ of a particle, while the corresponding covector $`p`$ describes the energy and momentum of the corresponding particle. Because of this physical importance, it is necessary to generalize the notion of covectors to the new space.
We would like to provide a generalization that preserves the physical meaning of the connection between vectors and covectors. The physical connection is probably best described in quantum mechanics. In quantum mechanics, due to Heisenberg’s uncertainty principle $`\mathrm{\Delta }x\mathrm{\Delta }p\mathrm{}`$, if we know the exact location of a particle (i.e., if $`\mathrm{\Delta }x=0`$), then we have no information about the momentum (i.e., $`\mathrm{\Delta }p=\mathrm{}`$), and vice versa, if we know the exact momentum ($`\mathrm{\Delta }p=0`$), then we have no information about the particle’s location. In other words, if we have a state with a definite momentum $`p`$, and we then shift the coordinates by a vector $`t`$, i.e., replace $`x`$ by $`x+t`$, the known state of the particle should not change.
In quantum mechanics, a state of the particle is described by a complex-valued function $`\psi (x)`$ called a wave function. The wave function itself is not directly observable, what we observe are probabilities $`|\psi |^2`$. So, if we multiply all the values of the wave-function by a complex number $`\phi `$ with $`|\phi |=1`$ (i.e., by a number of the type $`\mathrm{exp}(\mathrm{i}\alpha )`$, where $`\mathrm{i}=\sqrt{1}`$ and $`\alpha `$ is a real number), then all the probabilities remain the same – i.e., from the physical viewpoint, we will have exactly the same state. Thus, for every real number $`\alpha `$, the functions $`\psi (x)`$ and $`\mathrm{exp}(\mathrm{i}\alpha )\psi (x)`$ describe exactly the same state. When we say that the state $`\psi (x)`$ does not change after shift $`xx+t`$, we mean that the original function $`\psi (x)`$ and the function $`\psi (x+t)`$ that describe the shifted state describe the same state – i.e., $`\psi (x+t)=\phi (t)\psi (x)=\mathrm{exp}(\mathrm{i}\alpha (t))\psi (x)`$ for some complex number $`\phi (t)`$ or, equivalently, real number $`\alpha (t)`$ (which, generally speaking, depends on the shift $`t`$).
Since $`\mathrm{exp}(\mathrm{i}2\pi )=1`$, the value $`\alpha (t)`$ is only determined modulo $`2\pi `$. Thus, $`\alpha (t)`$ is a point on a circle rather than a real number.
For $`x=0`$, we get $`\psi (t)=\phi (t)\psi (0)`$, so modulo a multiplicative constant, shift-invariant states $`\psi (t)`$ are equal to the corresponding functions $`\phi (t)`$. So, to determine such states, we must describe all the corresponding functions $`\phi (t)`$.
When we shift by $`t=0`$, the function remains unchanged, i.e., $`\phi (0)=1`$ (equivalently, $`\alpha (0)=0`$).
If we first shift $`t`$ and then by $`s`$, then we get the same result as if we shift once by $`t+s`$. Hence, we have
$$\phi (s)(\phi (t)\psi (x))=\phi (t+s)\psi (x),$$
so $`\phi (t+s)=\phi (t)\phi (s)`$. So, from the physical viewpoint, a shift-invariant state $`\phi `$ is a mapping from $`R^4`$ to the unit circle $`S^1=\{\phi :|\phi |=1\}`$ that transform 0 into 1 and sum into sum. In mathematics, such a mapping is called a homomorphism from an Abelian additive group $`R^4`$ to $`S^1`$.
It is also physically reasonable to assume that the wave function is continuous – hence, that the homomorphism $`\phi `$ is continuous. Continuous homomorphisms from an Abelian group $`G`$ to a unit circle are called characters; the set of all such characters is also an Abelian group called dual (and denoted by $`G^{}`$). So, it is natural to associate covectors with elements of the dual group.
For $`R^4`$, this definition fits well with the more traditional one, because it is known that for $`R^4`$, the dual group is also $`R^4`$: every character has the form $`\mathrm{exp}(\mathrm{i}px)`$. For $`K=R^4\times S^1`$, we get a new definition:
Definition 2. A $`K`$-covector is a character of the group $`K`$, i.e., a continuous homomorphism from $`K`$ to $`S^1`$. By a sum of two covectors we mean the product of the corresponding homomorphisms.
The set of all $`K`$-covectors is thus a dual group $`K^{}`$ to $`K`$. It is known that elements of this dual group have the form $`\mathrm{exp}(\mathrm{i}px)`$, where $`p=(p_1,\mathrm{},p_4,p_5)`$, $`p_1,\mathrm{},p_4`$ can be any real numbers, and $`p_5`$ is an multiple of $`1/h`$. Thus, the group $`K^{}`$ of all $`K`$-covectors is isomorphic to $`R^4\times Z`$, where $`Z`$ is the additive group of all integers.
Comment. $`K`$-vectors are simply elements $`x=(x_1,\mathrm{},x_5)`$ of $`R^5`$, some of which are equivalent to each other: $`xx^{}`$ if $`x_5x_5^{}=kh`$ for some integer $`k`$. In other words, a $`K`$-vector can be viewed as a set
$$\{x^{}:x^{}x\}=\{(x_1,\mathrm{},x_4,x_5+kh\}.$$
A unit circle $`S^1`$ can also be described as simply the set $`R`$ of all real numbers with the equivalence relation $`\alpha \alpha ^{}`$ if and only if $`\alpha \alpha ^{}=k(2\pi )`$ – or, equivalently, as the class of sets $`\{\alpha +k(2\pi )\}`$.
In these terms, we can alternative describe $`K`$-covectors as linear mappings $`x=(x_1,\mathrm{},x_5)px=p_ix_i`$ from $`R^5`$ to $`R`$ that are consistent with the above structures, i.e., mapping for which $`xx^{}`$ implies $`pxpx^{}`$.
## 8 $`K`$-Tensors
To describe individual particles, it is usually sufficient to consider vectors (that describe their location) and covectors (that describe their momentum). However, to describe field theories such as Maxwell’s theory of electromagnetism or Einstein’s General Relativity theory, it is not sufficient to consider only vectors and covectors, we also need to consider tensors.
Specifically, for $`G=R^4`$, for every two integers $`p0`$ and $`q0`$, a tensor of valence $`(p,q)`$ can be defined as a multi-linear map $`G^p\times (G^{})^qR`$ – where multi-linear means that if we fix the values of all the variables but one, we get a linear mapping. Every such multi-linear mapping has the form
$$x^{i_1},\mathrm{},y^{i_p},z_{j_1},\mathrm{},u_{j_q}$$
$$\underset{i_1,\mathrm{},i_p,j_1,\mathrm{},j_q}{}t_{i_1\mathrm{}i_p}^{j_1\mathrm{}j_q}x^{i_1}\mathrm{}y^{i_p}z_{j_1}\mathrm{}u_{j_q}$$
for some components $`t_{i_1\mathrm{}i_p}^{j_1\mathrm{}j_q}`$. We thus naturally arrive at the following definition:
Definition 3. Let $`G_1,\mathrm{},G_m,G`$ be continuous Abelian groups. A mapping $`t:G_1\times \mathrm{}\times G_mG`$ is called $`Z`$-multilinear if for every $`i`$, if we fix the values of all the variables except $`i`$-th, we get a homomorphism.
Definition 4. Let $`p0`$ and $`q0`$. By a $`K`$-tensor of valence $`(p,q)`$, we mean a continuous $`Z`$-multilinear mapping $`t:K^p\times (K^{})^qS^1`$.
Comments. For $`R^4`$ instead of $`K`$, this definition coincides with the traditional one.
When $`K=R^4\times S^1`$, this definition is consistent with the previous ones: $`K`$-tensors of valence $`(0,1)`$ are $`K`$-covectors, and $`K`$-tensors of valence $`(1,0)`$ are $`K`$-vectors.
This definition can be reformulated as follows: a $`K`$-tensor is a multi-linear mapping that is consistent with the equivalence sets structure, i.e., for which $`xx^{},\mathrm{},yy^{}`$ implies that $`t(x,\mathrm{},y,z,\mathrm{},u)t(x^{},\mathrm{},y^{},z,\mathrm{},u).`$
Two multi-linear mappings $`t`$ and $`t^{}`$ describe the same $`K`$-tensor if $`t(x,\mathrm{},y,z,\mathrm{},u)t^{}(x,\mathrm{},y,z,\mathrm{},u)`$ for all $`x,\mathrm{},y,z,\mathrm{},u`$.
The following result describes all such mappings:
Proposition 1.
* Every $`K`$-tensor has the form
$$\mathrm{exp}\left(\mathrm{i}\underset{i_1,\mathrm{},i_p,j_1,\mathrm{},j_q}{}t_{i_1\mathrm{}i_p}^{j_1\mathrm{}j_q}x^{i_1}\mathrm{}y^{i_p}z_{j_1}\mathrm{}u_{j_q}\right)$$
* for some components $`t_{\mathrm{}}^{\mathrm{}}`$. In this representation, of all the components in which one of the lower indices is 5, only a component $`t_5^{5\mathrm{}5}`$ can be non-zero, and it can only take values $`2\pi h^{q1}k`$ for some integer $`k`$.
* Vice versa, if we have a set of components $`t_{\mathrm{}}^{\mathrm{}}`$ in which of all the components in which one of the lower indices is 5, only a component $`t_5^{5\mathrm{}5}`$ may be non-zero, its value is $`2\pi h^{q1}k`$ for some integer $`k`$, then the above formula defines a $`K`$-tensor.
* Two sets of components $`t_{\mathrm{}}^{\mathrm{}}`$ and $`s_{\mathrm{}}^{\mathrm{}}`$ define the same $`K`$-tensor if and only if all their components coincides with a possible exception of components $`t^{5\mathrm{}5}`$ and $`s^{5\mathrm{}5}`$ which may differ by $`2\pi h^qk`$ for an integer $`k`$.
Comment. For readers’ convenience, all the proofs are given in the Appendix.
## 9 Explaining the Condition $`g_{55}=\mathrm{const}`$ and the Fact that Metric Does Not Depend on $`x^5`$
For $`g_{ij}`$, Proposition 1 implies that $`g_{55}=g_{5i}=0`$. Thus, the above geometric formalism explains the first two physical assumptions that we wanted to explain: that $`g_{55}=0`$ and that the distance $`\mathrm{\Delta }s^2=\underset{i=1}{\overset{5}{}}\underset{j=1}{\overset{5}{}}g_{ij}\mathrm{\Delta }x_i\mathrm{\Delta }x_j`$ between the two points $`x`$ and $`x+\mathrm{\Delta }x`$ only depends on their first 4 coordinates.
## 10 Differential Formalism for $`K`$-Tensor Fields
Definition 5. By a $`K`$-tensor field $`f_{\mathrm{}}^{\mathrm{}}(x)`$ of valence $`(p,q)`$, we mean a mapping that assigns, to every point $`xK`$, a $`K`$-tensor $`f_{\mathrm{}}^{\mathrm{}}(x)`$ of this valence.
Most physics is described in the language of differential equations. It is known that for every tensor field $`t_{i_1\mathrm{}i_p}^{j_1\mathrm{}j_q}`$ of valence $`(p,q)`$, its gradient $`t_{i_1\mathrm{}i_p}^{j_1\mathrm{}j_q}/x^m`$ is also a tensor field – of valence $`(p+1,q)`$. This new field is called a gradient tensor field. It is therefore natural to give the following definition:
Definition 6. We say that a $`K`$-tensor field of valence $`(p,q)`$ is differentiable if the corresponding component tensor field is continuously differentiable, and its gradient field also defines a $`K`$-tensor field.
In other words, to differentiate a $`K`$-tensor field, we form the corresponding tensor field, differentiate it, and then interpret the result as a $`K`$-tensor field of valence $`(p+1,q)`$. When is this possible? The answer to this question is as follows:
Proposition 2. The $`K`$-tensor field is differentiable if and only if all its components $`t_{\mathrm{}}^{\mathrm{}}`$ do not depend on $`x^5`$, with the possible exception of the component $`t^{5\mathrm{}5}`$ which may have the form $`2\pi h^{q1}x^5+f(x_1,\mathrm{},x_4)`$.
## 11 Cylindricity Explained
As a result of Proposition 2, we conclude that for all the components $`t`$ (except for angular-valued ones), we have the cylindricity condition $`t_{\mathrm{}}^{\mathrm{}}/x^5=0`$. Thus, the cylindricity conditions is also explained by the geometric model.
## 12 Linear Coordinate Transformations
In the traditional affine geometry, in addition to shifts, we can also consider arbitrary linear coordinates transformations. In geometric terms, we can define these transformations as continuous automorphisms of the additive group $`K_0=R^4`$. We can define vectors and tensors as continuous homomorphisms $`T:K_0^p\times (K_0^{})^qS^1`$; in this case, e.g., standard formulas for transforming covectors (i.e., continuous homomorphisms $`g:K_0S^1`$) can be uniquely determined by the requirement that the value $`g(a)`$ be preserved under such a transformation, i.e., that $`g^{}(a^{})=g(a)`$. Similarly, the transformation law for tensors can be determined by the condition that
$$t^{}(a_1^{},\mathrm{},a_p^{},b_1^{},\mathrm{},b_q^{})=t(a_1,\mathrm{},a_p,b_1,\mathrm{},b_1).$$
$`(1)`$
Similarly, for $`K=R^4\times S^1`$, we can define a $`K`$-linear transformation as follows:
Definition 7. By a $`K`$-linear transformation, we mean a continuous automorphism of the additive group of $`K`$.
Proposition 3. Every $`K`$-linear transformation has the form
$$x^5\pm x^5+\underset{i=1}{\overset{4}{}}A_ix^i;x^i\underset{j=1}{\overset{4}{}}b_j^ix^j,(i4).$$
The corresponding tensor transformations can be defined by the condition (1). Once can see that in this case, the tensor components are transformed just like the normal tensor components. In particular, under the above $`K`$-linear transformation, a covector is transformed as follows:
$$x_5\pm x_5,x_i\underset{i=1}{\overset{4}{}}c_i^jx_jA_ix_5,$$
where $`c_i^j`$ is the matrix that is inverse to $`b_j^i`$.
## 13 General Coordinate Transformations
Definition 8. A smooth transformation $`s:KK`$ is admissible if and only if for each point $`xK`$, the corresponding tangent transformation
$$a^ia_{\mathrm{new}}^i=\underset{j=1}{\overset{5}{}}\frac{s^i}{x^j}_{|x}a^i$$
is a $`K`$-linear transformation.
Proposition 4. Every admissible transformation has the form
$$x^5\pm x^5+f(x^1,\mathrm{},x^4),x^if^i(x^1,\mathrm{},x^4).$$
Comment. We have already mentioned that functions on $`K=R^4\times S^1`$ are simply functions on $`R^5`$ which are periodic in $`x^5`$ with the period $`h`$. Also, a $`K`$-covector $`p`$ can be simply viewed as a covector for which the fifth component $`p_5`$ is an integer multiple of $`1/h`$. Thus,, e.g., a $`K`$-covector field on $`K`$ can be viewed as a covector field $`p(x)=(p_1(x),\mathrm{},p_5(x))`$ on $`R^5`$ that satisfies the following two properties:
* this field is periodic in $`x^5`$ with period $`p`$;
* for each $`x`$, the value $`p_5(x)`$ is an integer multiple of $`1/h`$.
It is therefore reasonable to define a general coordinate transformation of $`K`$ as a coordinate transformation of $`R^5`$ that preserves this property, i.e., under which a covector field that satisfies the properties (a) and (b) are transformed into a covector field that also satisfies these properties. One can see that this leads to the same class of general coordinate transformations.
## 14 Gauge Transformations Explained
According to Proposition 4, every admissible transformation is a composition of a 4D transformation and an additional gauge transformation $`x^5x^5+f(x^1,\mathrm{},x^4)`$ – exactly as described by Rumer.
## 15 Case of Curved Space-Time
In modern physics, space-time is a manifold, i.e., a topological space $`V`$ which is locally diffeomorphic to $`R^4`$. Since our basic model is not $`R^4`$, but $`K=R^4\times S^1`$, it is reasonable to define a $`K`$-manifold as a topological space that is locally diffeomorphic to $`K`$.
From the mathematical viewpoint, $`K`$ is $`R^5`$ factorized over the vector $`e=(0,\mathrm{},0,h)`$: i.e., $`ab`$ if and only if $`ab`$ is an integer multiple of $`e`$. Thus, a natural way to describe a $`K`$-manifold is to describe a standard 5D manifold in which we have a vector $`e(x)`$ in every tangent space – i.e., a manifold with an additional vector field.
In this case, every tangent space is isomorphic to $`K`$. Thus, a $`K`$-tensor field can be defined as a mapping that maps every point $`xV`$ into a $`K`$-tensor defined over the space $`K`$ which is tangent at $`x`$.
## 16 Auxiliary Result: Why There Is No Physically Useful Gravitational Analog of Hertz Potential
In electromagnetism, in addition to the electromagnetic file $`F_{ij}`$ and the potential $`A_i`$ from which this filed can be obtained by differentiation $`F_{ij}=A_i/x_jA_j/x_i`$, there is also a useful notion of a Hertz potential $`H^{ik}`$ for which $`A^i`$ can be obtained by differentiation $`A^i=\underset{k}{}H^{ik}/x^k`$.
In gravitation, the natural analogy of potentials $`A_i`$ is the gravity tensor filed $`g^{ij}`$. From the purely mathematical viewpoint, it is possible to introduce a gravitational analog of the Hertz potential: namely, there exists a tensor field $`\mathrm{\Pi }^{ijk}`$ for which
$$g^{ij}=\underset{k}{}\frac{\mathrm{\Pi }^{ijk}}{x^k};$$
$`(2)`$
see, e.g., \[Palchik 1969\]. However, in contrast to the electromagnetic case, this new potential does not seem to have any physical applications. Why?
Our explanation is simple: while (2) is impossible in the 4D case, it is no longer possible if we consider 5D $`K`$-tensor fields.
## Acknowledgments
The research was partially supported by NASA under cooperative agreement NCC5-209, by NSF grants EAR-0112968, EAR-0225670, and EIA-0321328, and by NIH grant 3T34GM008048-20S1.
The authors are thankful to all the participants of the special section of the October 1997 Montreal meeting of the American Mathematical Society, where physico-geometric aspects of this research were presented, for valuable comments; we are especially thankful to Prof. Abraham Ungar who organized this session, and to Yakov Eliashberg (Stanford) for important comments.
## Appendix: Proofs
### Proof of Proposition 1
Let us first prove that every $`K`$-tensor can be described by the desired formula.
Indeed, let $`t`$ be a $`K`$-tensor. Let us first consider the restriction of $`t`$ to $`K^p\times (R^4)^q`$. Since locally, $`K`$ coincides with $`R^5`$, this restriction is, locally, a multi-linear map from $`(R^5)^p\times (R^4)^q`$ to $`S^1`$. Since it is multi-linear, at 0, the value of this map is 1. In a small vicinity of 1, we can define a unique angle $`(1/\mathrm{i})\mathrm{ln}t`$. The resulting mapping is – locally – a multi-linear mapping, in the traditional sense of this term, from $`(R^5)^p\times (R^4)^q`$ to $`R`$. Hence, in this vicinity, $`\mathrm{ln}t=\mathrm{i}t_{\mathrm{}}^{\mathrm{}}x^{i_1}\mathrm{}`$ So, for the restriction of $`t`$ to $`K^p\times (R^4)^q`$, we get the desired formula.
Similarly, for $`K_0\stackrel{\mathrm{def}}{=}(R^5)^m\times R^4\times \mathrm{}\times R^4\mathrm{}\times e\times R^4\times \mathrm{}\times R^4`$, with $`r`$-th term replaced by $`e\stackrel{\mathrm{def}}{=}(0,0,0,0,h^1)`$, we conclude that the restriction of $`t`$ to $`K_0`$ has the form
$$\mathrm{exp}\left(\mathrm{i}t_{i_1\mathrm{}i_p}^{j_1\mathrm{}j_{r1}5j_{r+1}\mathrm{}j_q}x^{i_1}\mathrm{}y^{i_p}z_{j_1}\mathrm{}u_{j_q}\right)$$
for some values $`t_{i_1\mathrm{}i_p}^{j_1\mathrm{}j_{r1}5j_{r+1}\mathrm{}j_q}`$. Since the restriction of $`t`$ to the $`r`$-th copy of $`K^{}`$ is a homomorphism, this formula also holds for elements of $`(R^5)^m\times R^4\times \mathrm{}\times R^4\mathrm{}\times Z\times R^4\times \mathrm{}\times R^4`$,
Similar formulas hold for the subsets that can be obtained by replacing some of $`K^{}=R^4\times Z`$ with $`R^4`$ and some by $`Z`$. Since $`t`$ is a homomorphism w.r.t. each of its variables, we can represent each element $`p=(p_1,\mathrm{},p_4,p_5)K^{}`$ as a sum of $`p^{(4)}=(p_1,\mathrm{},p_4,0)R^4`$ and $`p^{(5)}=(0,\mathrm{},0,p_5)Z`$. For each of these two vectors, we have the desired formula; multiplying them, we get a similar formula for $`p`$. By using a similar decomposition w.r.t. other variables, we get the desired formula for all possible inputs from $`K^p\times (K^{})^q`$.
Let us now prove the desired properties of the components $`t_{\mathrm{}}^{\mathrm{}}`$. Since $`t`$ is defined on $`K^p\times (K^{})^q`$, replacing $`x^5`$ with $`x^5+h`$ should change the sum
$$\underset{i_1,i_2,\mathrm{},i_p,j_1,\mathrm{},j_q}{}t_{i_1i_2\mathrm{}i_p}^{j_1\mathrm{}j_q}x^{i_1}d^{i_2}\mathrm{}y^{i_p}z_{j_1}\mathrm{}u_{j_q}$$
by an integer multiple of $`2\pi `$. In other words, the difference between the new sum and old sum, i.e.,
$$h\underset{5,i_2,\mathrm{},i_p,j_1,\mathrm{},j_q}{}t_{5i_2\mathrm{}i_p}^{j_1\mathrm{}j_q}d^{i_2}\mathrm{}y^{i_p}z_{j_1}\mathrm{}u_{j_q}$$
must be a multiple of $`2\pi `$ for all $`d^{i_2},\mathrm{},y^{i_p}`$.
Let us first consider the case $`p>1`$. For $`d^{i_2}=\mathrm{}=y^{i_p}=0`$, the difference is equal to 0; this difference continuously depends on $`d^{i_2},\mathrm{},y^{i_p}`$, and it is only allowed a discrete set of values. Due to continuity, it cannot “jump” to values $`2\pi k`$ for $`k0`$, hence it is always equal to 0. So, the above polynomial is identically 0, hence all its coefficients $`t_{5i_2\mathrm{}i_p}^{j_1\mathrm{}j_q}`$ are identically 0.
Similarly, we can prove that $`t_5^{i_1\mathrm{}}=0`$ if $`i_15`$, so $`t_5^{5\mathrm{}5}`$ is indeed the only non-zero component of $`t_{\mathrm{}}^{\mathrm{}}`$ for which one of the lower indices is 5. For this component, the fact that $`ht_5^{5\mathrm{}5}p_5\mathrm{}p_5=2\pi k`$, where $`p_5=1/h`$, leads to the desired formula for $`t_5^{5\mathrm{}5}`$.
To complete the proof, let us assume that the two sets of coefficients $`t_{\mathrm{}}^{\mathrm{}}`$ and $`s_{\mathrm{}}^{\mathrm{}}`$ define the same $`K`$-tensor. This means that for their difference $`\delta _{\mathrm{}}^{\mathrm{}}`$, the sum
$$\underset{i_1,\mathrm{},i_p,j_1,\mathrm{},j_q}{}\delta _{i_1\mathrm{}i_p}^{j_1\mathrm{}j_q}x^{i_1}\mathrm{}y^{i_p}z_{j_1}\mathrm{}u_{j_q}$$
is an integer multiple of $`2\pi `$ for all $`x^{i_1},\mathrm{},y^{i_p}K`$ and $`z_{j_1},\mathrm{},u_{j_q}K^{}`$. If $`p>0`$, and one of the indices $`j_1,\mathrm{},j_q`$ is different from 5, then, as above, we can conclude that the sum is always 0,
So, all the corresponding coefficients $`\delta _{\mathrm{}}^{\mathrm{}}`$ are identically 0. The only possibly non-zero coefficient is $`\delta ^{5\mathrm{}5}`$. For this coefficient, the value $`\delta ^{5\mathrm{}5}p_5\mathrm{}p_5`$, with $`p_5=1/h`$, must be proportional to $`2\pi `$ – so $`\delta ^{5\mathrm{}5}(1/h)^p=2\pi k`$ for some integer $`k`$. Hence, the difference between $`s^{5\mathrm{}5}`$ and $`s^{5\mathrm{}5}`$ is indeed proportional to $`2\pi h^p`$. The proposition is proven.
### Proof of Proposition 2
According to Proposition 1, the only possibly non-zero component of a $`K`$-tensor with 5 as one of the lower indices is the component $`t_5^{5\mathrm{}5}`$. All the values $`t_{i_1\mathrm{}i_p}^{j_1\mathrm{}j_q}/x^5`$ contain 5 as one of the lower indices, so the only component for which this value can be different from 0 is the one with $`p=0`$ and $`i_1=\mathrm{}=i_p=5`$. For this component, $`t^{5\mathrm{}5}/x^5=2\pi h^{p1}k`$. Since the $`K`$-tensor field is continuously differentiable, this value cannot jump to a different value of $`k`$, so this derivative is constant. Integrating over $`x^5`$, we get the desired formula for the the dependence of this component on $`x^5`$ – as a linear function of $`x^5`$.
### Proof of Proposition 3
Since $`K`$ locally coincides with $`R^5`$, its continuous automorphisms locally coincide with continuous automorphisms $`R^5R^5`$, i.e., with linear transformations
$$x_{\mathrm{new}}^5=A_5x^5+\underset{i=1}{\overset{4}{}}A_ix^i;x_{\mathrm{new}}^i=B^ix^5+\underset{j=1}{\overset{4}{}}b_j^ix^j.$$
If $`y^5=x^5+h`$ and $`y^i=x^i`$ for all other $`i`$, then $`x`$ and $`y`$ define the exact same point in $`K`$. Therefore, the new values $`x_{\mathrm{new}}`$ and $`y_{\mathrm{new}}`$ must also define the same point, hence $`y_{\mathrm{new}}^i=x_{\mathrm{new}}^i`$ for $`i=1,\mathrm{},4`$ (hence $`B^i=0`$) and $`y_{\mathrm{new}}^5x_{\mathrm{new}}^5=`$ integer multiple of $`h`$ (hence $`A_5`$ is an integer).
Reversibility implies that $`A_5^1`$ should also be an integer, hence $`A_5=\pm 1.`$
### Proof of Proposition 4
The condition that the tangent transformation is $`K`$-linear means that $`s^5/x^5=\pm 1`$ (and due to continuity this does not depend on the point $`x`$, i.e., either it is everywhere equal to 1, or it is everywhere equal to $`1`$), and $`s^i/x^5=0`$ for $`i<4`$. Hence, $`s^5=\pm x_5+f(x^1,\mathrm{},x^4)`$ and $`s^i=f^i(x^1,\mathrm{},x^4)`$ for $`i<5`$.
### Proof of a the Statement About Hertz Potentials
As we have mentioned, it is possible that $`g^{55}/x^50`$. However, if the representation (2) was possible, then we would have
$$\frac{g^{55}}{x^5}=\frac{^2\mathrm{\Pi }^{555}}{(x^5)^2}+\underset{i=1}{\overset{4}{}}\frac{^2\mathrm{\Pi }^{55i}}{x^ix^5}.$$
However, according to our general result about components of $`K`$-tensors, all the terms in the right-hand side are 0s, so their sum cannot be equal to a non-zero value $`g^{55}/x^5`$.
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# Dressed coordinates: the path-integrals approach
## 1 Introduction
In recent works, the concept of dressed coordinates and dressed states have been introduced in the context of a harmonic oscillator linearly coupled to a massless scalar field . As emphasized in these references, the introduction of the dressed (or renormalized ) coordinate and state concepts is necessary in order to give physical consistence to the oscillator-field system as a toy model to describe an atom in interaction with the electromagnetic field, where the atom is roughly modeled by the harmonic oscillator. Also, as early stressed the introduction of dressed coordinates and states is twofold advantageous. From the physical view point, the dressed states behave as the one expected for the physically measurable states: excited atomic states are unstable whereas the atom in their ground state and no field quanta is stable. On the other hand, it allows exact computations for the probability amplitudes of the different radiation processes of the atom. Indeed, when the calculation is performed for weak coupling constant we obtain, for the spontaneous decay of the first excited state of the atom, the long know result: $`e^{\mathrm{\Gamma }t}`$. Furthermore, when applied to a confined atom, approximated by the harmonic oscillator, in a spherical cavity of sufficiently small diameter the method accounts for the experimentally observed inhibition of the decaying processes . Besides that, in Refs. the extension of the dressed coordinate and state concepts for non linear system have been addressed.
Nevertheless, in all previous works the approach used has been via the operatorial formalism of Quantum Mechanics. The aim of this paper is to develop a path-integral approach to the problem. We hope that the path integral approach will be more adequate in dealing with the problem of computing the reduced density matrix or to obtain the master equation for the atom as is the case in Caldeira-Legget type models . As in previous works, because its exact integrability, the dressed coordinates were introduced in the context of a harmonic oscillator coupled linearly to a massless scalar field.
The paper is organized as follows: In section 2 we introduce the model and compute the propagator by exact diagonalization. Section 3 is devoted to establish the dressed coordinates. In section 4 we compute the probabilities associated to some physical processes and extend the sum rules found in . Finally, in section 5 we give our concluding remarks. Through this paper we use natural units $`c=\mathrm{}=1`$.
## 2 The model and its exact diagonalization
We consider as a toy model of an atom-electromagnetic field system the system composed by a harmonic oscillator (the atom) coupled to a massless scalar field. By considering the dipole approximation and expanding in the field modes we get the following Hamiltonian
$$H=\frac{1}{2}\left(p_0^2+\omega _0^2q_0^2\right)+\frac{1}{2}\underset{k=1}{\overset{N}{}}\left(p_k^2+\omega _k^2q_k^2\right)q_0\underset{k=1}{\overset{N}{}}c_kq_k+\frac{1}{2}\underset{k=1}{\overset{N}{}}\frac{c_k^2}{\omega _k^2}q_0^2,$$
(1)
where $`q_0`$ is the oscillator coordinate and $`q_k`$ are the field modes with $`k=1,2,\mathrm{}`$; $`\omega _k=2\pi /L`$, $`c_k=\eta \omega _k`$, $`\eta =\sqrt{2g\mathrm{\Delta }\omega }`$, $`\mathrm{\Delta }\omega =\omega _{k+1}\omega _k=2\pi /L`$. With $`g`$ being a frequency dimensional coupling constant and $`L`$ the diameter of the sphere in which we confine the oscillator-field system. In Eq. (1) the limit $`N\mathrm{}`$ is to be understood. The last term in Eq. (1) can be seen as a frequency renormalization and, it guarantees a positive-defined Hamiltonian. Due to the Hamiltonian (1) is quadratic in the momenta and there are not constraints we can write the propagator for the system as being
$$K(\stackrel{}{q}_f,t;\stackrel{}{q}_i,0)=𝒟q_0\underset{k=1}{\overset{N}{}}𝒟q_k\mathrm{exp}\left(i_0^t𝑑tL\right),$$
(2)
where $`\stackrel{}{q}=(q_0,q_1,\mathrm{},q_{N1},q_N)^T`$ and $`L`$, the Lagrangian, is given by
$`L`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\dot{q}_0^2\overline{\omega }_0^2q_0^2\right)+{\displaystyle \frac{1}{2}}{\displaystyle \underset{k=1}{\overset{N}{}}}\left(\dot{q}_k^2\omega _k^2q_k^2+2q_0c_kq_k\right)`$ (3)
$`=`$ $`{\displaystyle \frac{1}{2}}\dot{\stackrel{}{q}}^T\dot{\stackrel{}{q}}{\displaystyle \frac{1}{2}}\stackrel{}{q}^T𝐀\stackrel{}{q},`$
where
$$\overline{\omega }_0^2=\omega _0^2+\underset{k=1}{\overset{N}{}}\frac{c_k^2}{\omega _k^2},$$
(4)
and $`𝐀`$ is a symmetric matrix whose components are given by
$$𝐀=\left(\begin{array}{ccccccc}\overline{\omega }_0^2& c_1& c_2& & & c_{N1}& c_N\\ c_1& \omega _1^2& & & & & \\ c_2& & \omega _2^2& & & & \\ \mathrm{}& & & \mathrm{}& & & \\ \mathrm{}& & & & \mathrm{}& & \\ c_{N1}& & & & & \omega _{N1}^2& \\ c_N& & & & & & \omega _N^2\end{array}\right).$$
(5)
To compute the propagator, Eq. (2), we introduce the coordinate transformation
$$\stackrel{}{q}=𝐓\stackrel{}{Q},q_\mu =\underset{r=0}{\overset{N}{}}t_\mu ^rQ_r,Q_r=\underset{\mu =0}{\overset{N}{}}t_\mu ^rq_\mu $$
(6)
where $`𝐓`$<sup>1</sup><sup>1</sup>1The components of matrix $`𝐓`$ and some of its properties are given in the Appendix. As T diagonalizes the matrix A, the coordinates $`Q_r`$ are known as *normal coordinates* is an orthogonal matrix that diagonalize $`𝐀`$,
$$𝐃=𝐓^T\mathrm{𝐀𝐓}=\text{diag}(\mathrm{\Omega }_0^2,\mathrm{\Omega }_1^2,\mathrm{\Omega }_2^2,\mathrm{},\mathrm{\Omega }_{N1}^2,\mathrm{\Omega }_N^2).$$
(7)
It is easy to show that the eigenvalues of $`𝐀`$, $`\mathrm{\Omega }_r`$ are obtained by solving the equation
$$\omega _0^2\mathrm{\Omega }^2=\eta ^2\underset{k=1}{\overset{N}{}}\frac{\mathrm{\Omega }^2}{\omega _k^2\mathrm{\Omega }^2}.$$
(8)
It has shown that such equation has definite positive frequencies $`\mathrm{\Omega }^2`$ as solutions . The solutions of the characteristic equation (8) were used to describe radiation process in small cavities in good agreement with the experiment. In this system is also used to describe a Brownian particle coupled to an ohmic environment.
Replacing Eq. (6) in Eq. (2) we get
$$K(\stackrel{}{q}_f,t;\stackrel{}{q}_i,0)=\underset{r=0}{\overset{N}{}}𝒟Q_r\mathrm{exp}\left[i_0^t𝑑t\left(\frac{1}{2}\dot{Q}_r^2\frac{1}{2}\mathrm{\Omega }_r^2Q_r^2\right)\right].$$
(9)
Note that in Eq. (9) the functional measure is maintained, this because $`det𝐓=1`$. By using the known result for the propagator of a harmonic oscillator we get for Eq. (9),
$$K(\stackrel{}{q}_f,t;\stackrel{}{q}_i,0)=\underset{\mu =0}{\overset{N}{}}\left(\frac{\mathrm{\Omega }_\mu }{i2\pi \mathrm{sin}\left(\mathrm{\Omega }_\mu t\right)}\right)^{\frac{1}{2}}\mathrm{exp}\left[\frac{i\mathrm{\Omega }_\mu }{2\mathrm{sin}\left(\mathrm{\Omega }_\mu t\right)}\left(\left[Q_{f\mu }^2+Q_{i\mu }^2\right]\mathrm{cos}\left(\mathrm{\Omega }_\mu t\right)2Q_{i\mu }Q_{f\mu }\right)\right].$$
(10)
The spectral function is defined as being
$$Y(t)=𝑑q_0𝑑q_1\mathrm{}𝑑q_NK(\stackrel{}{q},t;\stackrel{}{q},0)$$
(11)
which is easily computed expressing the integral in normal coordinates $`\left\{Q_r\right\}`$, thus we obtain
$$Y(t)=\underset{r=0}{\overset{N}{}}\left(i2\mathrm{sin}\left(\frac{\mathrm{\Omega }_rt}{2}\right)\right)^1=\underset{n_0,\mathrm{},n_N=0}{\overset{\mathrm{}}{}}e^{itE_{n_0,\mathrm{},n_N}}$$
(12)
with the energy spectrum being given by
$$E_{n_0,\mathrm{},n_N}=\underset{r=0}{\overset{N}{}}\mathrm{\Omega }_r\left(n_r+\frac{1}{2}\right).$$
(13)
In the $`\left\{Q_r\right\}`$ coordinates, the ground state wave function is computed from the propagator given by Eq. (10) by taking the limit $`ti\mathrm{}`$
$$K(\stackrel{}{Q}_f,ti\mathrm{};\stackrel{}{Q}_i,0)=\underset{r=0}{\overset{N}{}}\left(\frac{\mathrm{\Omega }_r}{\pi }\right)^{\frac{1}{2}}\mathrm{exp}(\frac{\mathrm{\Omega }_r}{2}[Q_{ir}^2+Q_{fr}^2]it\frac{\mathrm{\Omega }_r}{2}),$$
(14)
thus the wave function of the ground state is
$$\psi _{00\mathrm{}0}\left(\stackrel{}{Q}\right)=\underset{r=0}{\overset{N}{}}\left(\frac{\mathrm{\Omega }_r}{\pi }\right)^{\frac{1}{4}}\mathrm{exp}\left(\frac{\mathrm{\hspace{0.17em}1}}{2}\mathrm{\Omega }_rQ_r^2\right).$$
(15)
We can write the ground state eigenfunction in the original coordinates $`q_\mu `$, by using the third equation of (6)
$$\psi _{00\mathrm{}0}(\stackrel{}{q})=\left(\frac{\mathrm{\Omega }_0}{\pi }\right)^{\frac{1}{4}}\left(\frac{\mathrm{\Omega }_1}{\pi }\right)^{\frac{1}{4}}\mathrm{}\left(\frac{\mathrm{\Omega }_N}{\pi }\right)^{\frac{1}{4}}\mathrm{exp}\left(\frac{\mathrm{\hspace{0.17em}1}}{2}\underset{\mu ,\nu =0}{\overset{N}{}}\underset{r=0}{\overset{N}{}}\mathrm{\Omega }_rt_\mu ^rt_\nu ^rq_\mu q_\nu \right).$$
(16)
## 3 The dressed coordinates
We have observed the vacuum wave function in the $`\{q_\mu \}`$ and normal $`\left\{Q_r\right\}`$ coordinates and we make a question: Is it possible to find some new set of coordinates $`\{\overline{q}_\mu \}`$ relate to them what allow us to describe the oscillators with their non interacting characteristics? The answer is yes. In such one new set of coordinates the vacuum wave function must be given as
$$\psi _{\mathrm{00..0}}(\overline{q}_0,\mathrm{},\overline{q}_N)=cte.\mathrm{exp}\left(\frac{1}{2}\underset{\mu =0}{\overset{N}{}}\omega _\mu \left(\overline{q}_\mu \right)^2\right)$$
(17)
and we named the set of coordinates $`\left\{\overline{q}_\mu \right\}`$ as *dressed coordinates*<sup>2</sup><sup>2</sup>2The dressed coordinates are a new type of collective coordinates, which are defined from the normal coordinates, that allows a correct description for the absorption and radiation phenomena for a given physical system. which describe the oscillators as being non interacting. In a physical situation we can imagine the atom interacting with a bath of field modes, i.e., an electromagnetic field. The experience tell us that the atom does not change his energy spectrum (the energy spectrum when it is isolated) and only transitions between the energy levels are observed (absorption and emission process).
Thus we would have for the vacuum state wave functions the following relation
$$\mathrm{exp}\left(\frac{1}{2}\underset{r=0}{\overset{N}{}}\mathrm{\Omega }_r\left(Q_r\right)^2\right)\mathrm{exp}\left(\frac{1}{2}\underset{\mu =0}{\overset{N}{}}\omega _\mu \left(\overline{q}_\mu \right)^2\right)$$
(18)
First we will look for the matrix transformation between the normal $`\left\{Q_\mu \right\}`$ and dressed coordinates $`\left\{\overline{q}_\mu \right\}`$, thus we set
$$\stackrel{}{\overline{q}}=𝐌\stackrel{}{Q},\overline{q}_\mu =\underset{r=0}{\overset{N}{}}M_\mu ^rQ_r$$
(19)
the quadratic form (18) must be invariant under the linear transformation (19), then, to preserve the quadric form we set
$$\underset{\mu =0}{\overset{N}{}}\omega _\mu M_\mu ^rM_\mu ^s=\mathrm{\Omega }_r\delta _{rs}$$
(20)
and to achieve the matrix $`𝐌`$ we use the orthonormal matrix $`𝐓`$ such that
$$M_\mu ^r=\sqrt{\frac{\mathrm{\Omega }_r}{\omega _\mu }}t_\mu ^r$$
(21)
we can show that it satisfies the equation (20) by using the orthonormality condition of the matrix $`𝐓`$. The determinant of the matrix $`𝐌`$ is shown to be 1,
$$det(𝐌)=\sqrt{\frac{\mathrm{\Omega }_0\mathrm{\Omega }_1\mathrm{}\mathrm{\Omega }_N}{\omega _0\omega _1\mathrm{}\omega _N}}det(𝐓)=1$$
thus the transformation (19) preserve the path-integral measure defining the propagator (2) of the system. And the inverse transformation $`\left\{Q_s\overline{q}_\mu \right\}`$ is easily shown to be
$$Q_s=\underset{\mu =0}{\overset{N}{}}\sqrt{\frac{\omega _\mu }{\mathrm{\Omega }_s}}t_\mu ^s\overline{q}_\mu $$
(22)
Then, we come back to the functional integral defining the propagator (9) in terms of the normal coordinates and using the transformation (22) we can write
$$K(\stackrel{}{q}_f,t;\stackrel{}{q}_i,0)=\left(\underset{\mu =0}{\overset{N}{}}𝒟\overline{q}_\mu \right)\mathrm{exp}\left\{i_0^t𝑑t\frac{1}{2}\underset{\mu ,\nu =0}{\overset{N}{}}\left[\frac{}{}Z_{\mu \nu }\dot{\overline{q}}_\mu \dot{\overline{q}}_\nu C_{\mu \nu }\omega _\mu \omega _\nu \overline{q}_\mu \overline{q}_\nu \right]\right\}$$
(23)
from which we can see the Lagrangian in terms of dressed coordinates is given by
$$L_d=\frac{1}{2}\underset{\mu ,\nu =0}{\overset{N}{}}\left[\frac{}{}Z_{\mu \nu }\dot{\overline{q}}_\mu \dot{\overline{q}}_\nu C_{\mu \nu }\omega _\mu \omega _\nu \overline{q}_\mu \overline{q}_\nu \right]$$
(24)
where we have defined the matrices $`Z_{\mu \nu }`$ and $`C_{\mu \nu }`$
$$Z_{\mu \nu }=\sqrt{\omega _\mu \omega _\nu }\underset{s=0}{\overset{N}{}}\frac{t_\mu ^st_\nu ^s}{\mathrm{\Omega }_s},C_{\mu \nu }=\frac{1}{\sqrt{\omega _\mu \omega _\nu }}\underset{s=0}{\overset{N}{}}\mathrm{\Omega }_st_\mu ^st_\nu ^s$$
(25)
such that
$$Z_{\mu \beta }C_{\beta \nu }=\delta _{\mu \nu }$$
(26)
They play the role of renormalization constants and the coordinates $`\left\{\overline{q}_\mu \right\}`$ are the renormalized coordinates such as it happened in a renormalized field theory.
To construct the dressed Hamiltonian we first define the dressed momentum $`\overline{p}_\mu `$ canonically conjugate to the dressed coordinate $`\overline{q}_\mu `$, thus
$$\overline{p}_\mu =\frac{L_d}{\dot{\overline{q}}_\mu }=\underset{\nu =0}{\overset{N}{}}Z_{\mu \nu }\dot{\overline{q}}_\nu $$
(27)
from which we get
$$\dot{\overline{q}}_\mu =\underset{\nu =0}{\overset{N}{}}C_{\mu \nu }\overline{p}_\nu $$
(28)
Thus, the dressed Hamiltonian $`H_d`$ is computed to be
$$H_d=\underset{\mu ,\nu =0}{\overset{N}{}}\frac{1}{2}C_{\mu \nu }\left(\frac{}{}\overline{p}_\mu \overline{p}_\nu +\omega _\mu \omega _\nu \overline{q}_\mu \overline{q}_\nu \right)$$
(29)
From (10) we write the propagator in dressed coordinates
$`K(\overline{q}_f,T;\overline{q}_i,0)`$ $`=`$ $`{\displaystyle \underset{r=0}{\overset{N}{}}}\left({\displaystyle \frac{\omega _r}{i2\pi \mathrm{sin}\left(\mathrm{\Omega }_rT\right)}}\right)^{\frac{1}{2}}`$
$`\times \mathrm{exp}\left({\displaystyle \frac{i}{2}}{\displaystyle \underset{s=0}{\overset{N}{}}}{\displaystyle \underset{\mu ,\nu =0}{\overset{N}{}}}{\displaystyle \frac{\sqrt{\omega _\mu \omega _\nu }t_\mu ^st_\nu ^s}{\mathrm{sin}\left(\mathrm{\Omega }_sT\right)}}\left[{\displaystyle \frac{}{}}\left(\overline{q}_{f\mu }\overline{q}_{f\nu }+\overline{q}_{i\mu }\overline{q}_{i\nu }\right)\mathrm{cos}\left(\mathrm{\Omega }_sT\right)2\overline{q}_{f\mu }\overline{q}_{i\nu }\right]\right)`$
It is simple to show that the spectral function computed in dressed coordinates is the same computed in the equation (12), thus, the energy spectrum (13) remains invariant, as it was expected.
The ground state reads as
$$\psi _{\mathrm{00..0}}\left(\overline{q}\right)=\left(\frac{\omega _0}{\pi }\right)^{1/4}\mathrm{}\left(\frac{\omega _N}{\pi }\right)^{1/4}\mathrm{exp}\left(\frac{1}{2}\underset{\alpha =0}{\overset{N}{}}\omega _\alpha \left(\overline{q}_\alpha \right)^2\right)$$
(31)
The Hamiltonian operator in dressed coordinates is expressed as
$$H\left(\overline{q}\right)=\underset{\mu ,\nu =0}{\overset{N}{}}C_{\mu \nu }\left(\frac{1}{2}\frac{}{\overline{q}_\nu }\frac{}{\overline{q}_\mu }+\frac{1}{2}\omega _\mu \omega _\nu \overline{q}_\mu \overline{q}_\nu \right).$$
(32)
with the coefficients $`C_{\mu \nu }`$ are given in (25).
## 4 Computing the transition probabilities
In this section we show what to compute the transition amplitudes of the system using the exact dressed propagator (3), thus, we are interested in the following quantities,
$$𝒜_{m_0m_1\mathrm{}m_N}^{n_0n_1\mathrm{}n_N}(t)={}_{d}{}^{}n_0,n_1,\mathrm{},n_N|e^{iHt}|m_0,m_1,\mathrm{},m_N_{d}^{},$$
(33)
that represents the probability amplitude of the system, initially prepared in the state $`|m_0,m_1,\mathrm{},m_N`$, to be found at time $`t`$ in the state $`|n_0,n_1,\mathrm{},n_N`$.
Eq. (33) can be written in terms of the propagator as
$$𝒜_{m_0m_1\mathrm{}m_N}^{n_0n_1\mathrm{}n_N}(t)=d\chi d\xi {}_{d}{}^{}n_0,n_1,\mathrm{},n_N|\chi K(\chi ,t;\xi ,0)\xi |m_0,m_1,\mathrm{},m_N_d,$$
(34)
where
$$\xi |m_0,m_1,\mathrm{},m_N_d=\psi _{m_0m_1\mathrm{}m_N}(\xi ^{}(\xi )).$$
(35)
Using Eq. (35) we can write Eq. (34) as
$$𝒜_{m_0m_1\mathrm{}m_N}^{n_0n_1\mathrm{}n_N}(t)=𝑑\chi 𝑑\xi \psi _{n_0n_1\mathrm{}n_N}(\chi ^{}(\chi ))K(\chi ,t;\xi ,0)\psi _{m_0m_1\mathrm{}m_N}(\xi ^{}(\xi ))$$
(36)
First we compute $`𝒜_{0\mathrm{}0m_\mu 0\mathrm{}0}^{0\mathrm{}0n_\nu 0\mathrm{}0}(t)`$. Substituting $`\psi _{0\mathrm{}0n_\nu 0\mathrm{}0}(\chi ^{}(\chi ))`$, $`K(\chi ,t;\xi ,0)`$ and $`\psi _{0\mathrm{}0m_\mu 0\mathrm{}0}(\xi ^{}(\xi ))`$ we have
$$𝒜_{0\mathrm{}0m_\mu 0\mathrm{}0}^{0\mathrm{}0n_\nu 0\mathrm{}0}(t)=\left[\underset{r=0}{\overset{N}{}}\frac{1}{\pi \sqrt{2i\mathrm{sin}(\mathrm{\Omega }_rt)}}\right]𝑑\chi \frac{H_{n_\nu }\left(_{r=0}^Nt_\nu ^r\chi _r\right)}{\sqrt{2^{n_\nu }n_\nu !}}\mathrm{exp}\left(\frac{i}{2}\underset{r=0}{\overset{N}{}}\frac{e^{i\mathrm{\Omega }_rt}}{\mathrm{sin}(\mathrm{\Omega }_rt)}\chi _r^2\right)F_\mu (\chi ,t)$$
(37)
where
$`F_\mu (\chi ,t)`$ $`=`$ $`{\displaystyle 𝑑\xi \frac{H_{m_\mu }\left(_{s=0}^Nt_\mu ^s\xi _s\right)}{\sqrt{2^{m_\mu }m_\mu !}}\mathrm{exp}\left[\frac{i}{2}\underset{s=0}{\overset{N}{}}\left(\frac{e^{i\mathrm{\Omega }_st}}{\mathrm{sin}(\mathrm{\Omega }_st)}\xi _s^2\frac{2\chi _s}{\mathrm{sin}(\mathrm{\Omega }_st)}\xi _s\right)\right]}`$ (38)
$`=`$ $`{\displaystyle \frac{H_{m_\mu }\left(_{s=0}^Nt_\mu ^si\mathrm{sin}(\mathrm{\Omega }_st)\frac{}{\chi _s}\right)}{\sqrt{2^{m_\mu }m_\mu !}}}{\displaystyle 𝑑\xi \mathrm{exp}\left[\frac{i}{2}\underset{s=0}{\overset{N}{}}\left(\frac{e^{i\mathrm{\Omega }_st}}{\mathrm{sin}(\mathrm{\Omega }_st)}\xi _s^2\frac{2\chi _s}{\mathrm{sin}(\mathrm{\Omega }_st)}\xi _s\right)\right]}.`$
Performing the Gaussian integrals in Eq. (38) we obtain
$$F_\mu (\chi ,t)=\left[\underset{s=0}{\overset{N}{}}\frac{\sqrt{2\pi i\mathrm{sin}(\mathrm{\Omega }_st)}}{e^{\frac{i}{2}\mathrm{\Omega }_st}}\right]\frac{H_{m_\mu }\left(_{s=0}^Nt_\mu ^si\mathrm{sin}(\mathrm{\Omega }_st)\frac{}{\chi _s}\right)}{\sqrt{2^{m_\mu }m_\mu !}}\mathrm{exp}\left(\frac{i}{2}\underset{r=0}{\overset{N}{}}\frac{e^{i\mathrm{\Omega }_rt}}{\mathrm{sin}(\mathrm{\Omega }_rt)}\chi _r^2\right).$$
(39)
Using the identity
$$H_n(\underset{r=0}{\overset{N}{}}t_\mu ^rX_r)=n!\underset{+l=n}{}\frac{(t_\mu ^0)^{l_0}}{l_0!}\frac{(t_\mu ^1)^{l_1}}{l_1!}\mathrm{}\frac{(t_\mu ^N)^{l_N}}{l_N!}H_{l_0}(X_0)H_{l_1}(X_1)\mathrm{}H_{l_N}(X_N)$$
(40)
(that holds because $`\{t_\mu ^r\}`$ is an orthogonal matrix) where $`+l=(l_0+l_1+\mathrm{}+l_N)`$ and replacing Eq. (39) in Eq. (37) we get
$$𝒜_{0\mathrm{}0m_\mu 0\mathrm{}0}^{0\mathrm{}0n_\nu 0\mathrm{}0}(t)=\pi ^{(N+1)/2}e^{\frac{i}{2}_{r=0}^N\mathrm{\Omega }_rt}\sqrt{\frac{m_\mu !n_\nu !}{2^{m_\mu +n_\nu }}}\underset{+l=n_\nu }{}\underset{+s=m_\mu }{}\frac{(t_\mu ^0)^{l_0+s_0}}{l_0!s_0!}\frac{(t_\mu ^1)^{l_1+s_1}}{l_1!s_1!}\mathrm{}\frac{(t_\mu ^N)^{l_N+s_N}}{l_N!s_N!}I_{l_0s_0}I_{l_1s_1}\mathrm{}I_{l_Ns_N}$$
(41)
where
$`I_{l_rs_r}`$ $`=`$ $`{\displaystyle 𝑑\chi _r\mathrm{exp}\left(\frac{i}{2}\frac{e^{i\mathrm{\Omega }_rt}}{\mathrm{sin}(\mathrm{\Omega }_rt)}\chi _r^2\right)H_{l_r}(\chi _r)H_{s_r}\left(i\mathrm{sin}(\mathrm{\Omega }_rt)\frac{}{\chi _r}\right)\mathrm{exp}\left(\frac{i}{2}\frac{e^{i\mathrm{\Omega }_rt}}{\mathrm{sin}(\mathrm{\Omega }_rt)}\chi _r^2\right)}`$ (42)
$`=`$ $`{\displaystyle 𝑑\chi _re^{\chi _r^2}H_{l_r}(\chi _r)\left[\mathrm{exp}\left(\frac{i}{2}\frac{e^{i\mathrm{\Omega }_rt}}{\mathrm{sin}(\mathrm{\Omega }_rt)}\chi _r^2\right)H_{s_r}\left(i\mathrm{sin}(\mathrm{\Omega }_rt)\frac{}{\chi _r}\right)\mathrm{exp}\left(\frac{i}{2}\frac{e^{i\mathrm{\Omega }_rt}}{\mathrm{sin}(\mathrm{\Omega }_rt)}\chi _r^2\right)\right]}.`$
If instead of integrating over coordinates $`\xi `$ in Eq. (37) we first integrate over coordinates $`\chi `$ we would get an expression similar to the one given in Eq. (41) but with $`I_{l_rs_r}`$ replaced with $`I_{l_rs_r}^{}`$:
$$I_{l_rs_r}^{}=𝑑\xi _re^{\chi _r^2}H_{s_r}(\xi _r)\left[\mathrm{exp}\left(\frac{i}{2}\frac{e^{i\mathrm{\Omega }_rt}}{\mathrm{sin}(\mathrm{\Omega }_rt)}\xi _r^2\right)H_{l_r}\left(i\mathrm{sin}(\mathrm{\Omega }_rt)\frac{}{\xi _r}\right)\mathrm{exp}\left(\frac{i}{2}\frac{e^{i\mathrm{\Omega }_rt}}{\mathrm{sin}(\mathrm{\Omega }_rt)}\xi _r^2\right)\right].$$
(43)
Then, since the final result must not depend of the order in which we perform the integrations we must have $`I_{l_rs_r}=I_{l_rs_r}^{}`$, and from Eqs. (42) and (43) we conclude that $`I_{l_rs_r}=I_{s_rl_r}`$.
To perform the integral given in Eq. (42) we have to use the following theorem
$$\mathrm{if}k<n𝑑xe^{x^2}H_n(x)x^k=0.$$
(44)
Note that the expression in brackets in Eq. (42) is a polynomial of degree $`s_r`$ in $`\xi _r`$. Now, if $`l_r>s_r`$, then by using theorem (44), we get $`I_{l_rs_r}=0`$. Because $`I_{l_rs_r}=I_{s_rl_r}`$ we also get a vanishing result for $`l_r<s_r`$. Then, the only non vanishing result is obtained for $`l_r=s_r`$. Using again theorem (44) we note that the only non vanishing term of the polynomial in brackets is the one of highest power. Since the highest power of $`H_n(x)`$ is given by $`2^nx^n`$ we have for Eq. (42)
$`I_{l_rs_r}`$ $`=`$ $`[2i\mathrm{sin}(\mathrm{\Omega }_rt)]^{s_r}{\displaystyle 𝑑\chi _re^{\chi _r^2}H_{l_r}(\chi _r)\left[\mathrm{exp}\left(\frac{i}{2}\frac{e^{i\mathrm{\Omega }_rt}}{\mathrm{sin}(\mathrm{\Omega }_rt)}\chi _r^2\right)\frac{^{s_r}}{\chi _r^{s_r}}\mathrm{exp}\left(\frac{i}{2}\frac{e^{i\mathrm{\Omega }_rt}}{\mathrm{sin}(\mathrm{\Omega }_rt)}\chi _r^2\right)\right]}`$ (45)
$`=`$ $`e^{is_r\mathrm{\Omega }_rt}{\displaystyle 𝑑\chi _re^{\chi _r^2}H_{l_r}(\chi _r)(2)^{s_r}\chi _r^{s_r}}`$
$`=`$ $`e^{is_r\mathrm{\Omega }_rt}{\displaystyle 𝑑\chi _re^{\chi _r^2}H_{l_r}(\chi _r)H_{s_r}(\chi _r)}`$
$`=`$ $`\sqrt{\pi }e^{is_r\mathrm{\Omega }_rt}2^{s_r}s_r!\delta _{l_rs_r}.`$
Using Eq. (45) in Eq. (41) we get
$`𝒜_{0\mathrm{}0m_\mu 0\mathrm{}0}^{0\mathrm{}0n_\nu 0\mathrm{}0}(t)`$ $`=`$ $`e^{\frac{i}{2}_{r=0}^N\mathrm{\Omega }_rt}\sqrt{{\displaystyle \frac{m_\mu !n_\nu !}{2^{m_\mu +n_\nu }}}}{\displaystyle \underset{+l=m_\mu =n_\nu }{}}{\displaystyle \frac{((t_\mu ^0)^2e^{i\mathrm{\Omega }_0})^{l_0}}{l_0!}}{\displaystyle \frac{((t_\mu ^1)^2e^{i\mathrm{\Omega }_1})^{l_1}}{l_1!}}\mathrm{}{\displaystyle \frac{((t_\mu ^N)^2e^{i\mathrm{\Omega }_N})^{l_N}}{l_N!}}2^{l_0+l_1+\mathrm{}+l_N}`$ (46)
$`=`$ $`e^{\frac{i}{2}_{r=0}^N\mathrm{\Omega }_rt}\delta _{mn}\left({\displaystyle \underset{r=0}{\overset{N}{}}}t_\mu ^rt_\nu ^re^{i\mathrm{\Omega }_rt}\right)^n,`$
where in passing to the last line we have used the identity
$$\left(\underset{r=0}{\overset{N}{}}X_r\right)^n=n!\underset{+l=n}{}\frac{X_0^{l_0}}{l_0!}\frac{X_1^{l_1}}{l_1!}\mathrm{}\frac{X_N^{l_N}}{l_N!}.$$
(47)
In terms of
$$f_{\mu \nu }(t)=\underset{r=0}{\overset{N}{}}t_\mu ^rt_\nu ^re^{i\mathrm{\Omega }_rt},$$
(48)
Eq. (46) can be written as
$$𝒜_{0\mathrm{}0m_\mu 0\mathrm{}0}^{0\mathrm{}0n_\nu 0\mathrm{}0}(t)=e^{\frac{i}{2}_{r=0}^N\mathrm{\Omega }_rt}\delta _{mn}\left[f_{\mu \nu }(t)\right]^n.$$
(49)
It is straightforward to establish the following identity:
$$\underset{\mu =0}{\overset{N}{}}|f_{\nu \mu }(t)|^2=1$$
(50)
The proof of the above identity follows trivially by using the orthonormality property of the matrix $`\{t_\mu ^r\}`$. Writing Eq. (50) for indexes $`0`$ and $`k`$ we have
$$|f_{00}(t)|^2+\underset{k=1}{\overset{N}{}}|f_{0k}(t)|^2=1,$$
(51)
$$|f_{k_10}(t)|^2+\underset{k_2=1}{\overset{N}{}}|f_{k_1k_2}(t)|^2=1.$$
(52)
The physical interpretation for the equations above is given as it follows. Let the initial state of the system given by $`|n,1_{k_1},1_{k_2}\mathrm{}`$, the atom in the $`n`$-th excited level and field quanta of frequencies $`\omega _{k_1}`$, $`\omega _{k_2}`$, etc. The probability of this initial states to be found in a measurement performed at time $`t`$ in the state $`|m,1_{k_1^{}},1_{k_2^{}}\mathrm{}`$ is denoted by $`𝒫_{n;1_{k_1}1_{k_2}\mathrm{}}^{m;1_{k_1^{}}1_{k_2^{}}\mathrm{}}(t)`$. We know that $`𝒫_{1;0}^{1;0}(t)=|f_{00}(t)|^2`$ is the probability of the oscillator to remain in the first excited level and $`P_{1;0}^{0;1_k}(t)=|f_{0k}(t)|^2`$ is the probability of the oscillator to decay from the first excited level to the ground state by emission of a field quanta of frequency $`\omega _k`$. Obviously in this case we have
$$𝒫_{1;0}^{1;0}(t)+\underset{k}{}P_{1;0}^{0;1_k}(t)=1$$
(53)
that is nothing but Eq. (51). Also Eq. (52) can be written as
$$𝒫_{0;k_1}^{1;0}(t)+\underset{k_2}{}𝒫_{0;k_1}^{0;k_2}(t)=1,$$
(54)
where
$$𝒫_{0;k_1}^{1;0}(t)=|f_{k_10}(t)|^2$$
(55)
and
$$𝒫_{0;k_1}^{0;k_2}(t)=|f_{k_1k_2}(t)|^2.$$
(56)
With these identifications the physical meaning of Eq. (51) is clear: if initially we have a photon of frequency $`\omega _{k_1}`$ and the oscillator is in its ground state, then at time $`t`$, either the oscillator can go to its first excited level by absorption of the initial photon or can remain in its ground state scattering the initial photon to other photon of arbitrary frequency.
Note that in establishing the identities (53) and (54) it is used only the orthogonality property of the matrix $`\{t_\mu ^r\}`$. Then, it is natural to ask whether it is possible to compute other probabilities without doing a direct computation as performed in last section. The answer is yes. For example, if initially the oscillator is in its second excited level and there are no photons, at time $`t`$ it can happening that the oscillator continues in their second excited level, it can go to their first excited level by emission of photon of arbitrary frequency $`\omega _{k_1}`$ or it can decay to their ground state by emission of two photons of arbitrary frequencies $`\omega _{k_1}`$ and $`\omega _{k_2}`$. The respective probabilities are denoted by $`𝒫_{2;0}^{2;0}(t)`$, $`𝒫_{2;0}^{1;1_{k_1}}(t)`$ and $`𝒫_{2;0}^{0;1_{k_1}1_{k_2}}(t)`$. Obviously we must have
$$𝒫_{2;0}^{2;0}(t)+\underset{k_1}{}𝒫_{2;0}^{1;1_{k_1}}(t)+\underset{k_1k_2}{}𝒫_{2;0}^{0;1_{k_1}1_{k_2}}(t)=1.$$
(57)
Taking the square of Eq. (53) we find
$$\left(𝒫_{1;0}^{1;0}(t)\right)^2+2𝒫_{1;0}^{1;0}(t)\underset{k_1}{}P_{1;0}^{0;1_{k_1}}(t)+\underset{k_1k_2}{}P_{1;0}^{0;1_{k_1}}(t)P_{1;0}^{0;1_{k_2}}(t)=1.$$
(58)
Identifying Eqs. (57) and (58) we obtain
$`𝒫_{2;0}^{2;0}(t)`$ $`=`$ $`\left(𝒫_{1;0}^{1;0}(t)\right)^2`$ (59)
$`=`$ $`|f_{00}(t)|^4`$
$`𝒫_{2;0}^{1;1_{k_1}}(t)`$ $`=`$ $`2𝒫_{1;0}^{1;0}(t)P_{1;0}^{0;1_{k_1}}(t)`$ (60)
$`=`$ $`2|f_{00}(t)f_{0k_1}(t)|^2`$
and
$`𝒫_{2;0}^{0;1_{k_1}1_{k_2}}(t)`$ $`=`$ $`P_{1;0}^{0;1_{k_1}}(t)P_{1;0}^{0;1_{k_2}}(t)`$ (61)
$`=`$ $`|f_{0k_1}(t)f_{0k_2}(t)|^2.`$
As a second example we consider the oscillator is in its first excited state and there is one photon of frequency $`\omega _{k_1}`$. At time $`t`$ it can happen that: the oscillator go to its second excited level by absorbing the initial photon; or the oscillator remains in its first excited state and the initial photon is scattered to other photon of arbitrary frequency $`\omega _{k_2}`$; or maybe the oscillator can be decay to its ground state by emission of a photon of arbitrary frequency $`\omega _{k_2}`$ and the initial photon is scattered to other photon of frequency $`\omega _{k_3}`$. The respective probabilities are denoted by $`𝒫_{1;1_{k_1}}^{2;0}(t)`$, $`𝒫_{1;1_{k_1}}^{1;1_{k_2}}(t)`$ and $`𝒫_{1;1_{k_1}}^{0;1_{k_2}1_{k_3}}(t)`$. Then, we must have
$$𝒫_{1;1_{k_1}}^{2;0}(t)+\underset{k_2}{}𝒫_{1;1_{k_1}}^{1;1_{k_2}}(t)+\underset{k_2k_3}{}𝒫_{1;1_{k_1}}^{0;1_{k_2}1_{k_3}}(t).$$
(62)
Taking Eq. (53) times Eq. (54) we have
$$𝒫_{1;0}^{1;0}(t)𝒫_{0;1_{k_1}}^{1;0}(t)+\underset{k_2}{}\left(𝒫_{1;0}^{1;0}(t)𝒫_{0;1_{k_1}}^{0;1_{k_2}}(t)+𝒫_{0;1_{k_1}}^{1;0}(t)𝒫_{1;0}^{0;1_{k_2}}(t)\right)+\underset{k_2k_3}{}𝒫_{1;0}^{0;1_{k_2}}(t)𝒫_{0;1_{k_1}}^{0;1_{k_3}}(t)=1.$$
(63)
From Eqs. (62) and (63) we have
$`𝒫_{1;1_{k_1}}^{2;0}(t)`$ $`=`$ $`𝒫_{1;0}^{1;0}(t)𝒫_{0;1_{k_1}}^{1;0}(t)`$ (64)
$`=`$ $`|f_{00}(t)f_{0k_1}(t)|^2,`$
$`𝒫_{1;1_{k_1}}^{1;1_{k_2}}(t)`$ $`=`$ $`𝒫_{1;0}^{1;0}(t)𝒫_{0;1_{k_1}}^{0;1_{k_2}}(t)+𝒫_{0;1_{k_1}}^{1;0}(t)𝒫_{1;0}^{0;1_{k_2}}(t)`$ (65)
$`=`$ $`|f_{00}(t)f_{k_1k_2}(t)|^2+|f_{0k_1}(t)f_{0k_2}(t)|^2`$
and
$`𝒫_{1;1_{k_1}}^{0;1_{k_2}1_{k_3}}(t)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(𝒫_{1;0}^{0;1_{k_2}}(t)𝒫_{0;1_{k_1}}^{0;1_{k_3}}(t)+𝒫_{1;0}^{0;1_{k_3}}(t)𝒫_{0;1_{k_1}}^{0;1_{k_2}}(t)\right)`$ (66)
$`=`$ $`{\displaystyle \frac{1}{2}}\left(|f_{0k_2}(t)f_{k_1k_3}|^2+|f_{0k_3}(t)f_{k_1k_2}|^2\right).`$
And so we can give all the probabilities associated to any decay or absorption process placing in the system.
## 5 Conclusions
In the present work is shown using the path-integral formalism that the dressed coordinates appear as a coordinate transformation preserving the quadric form that defines the ground state wave function of the system guaranteeing the vacuum stability and, they also leave invariant the functional measure of the path-integral. Within the Hamiltonian formalism it can be shown that such linear transformation, which defines the dressed coordinates, also leaves invariant the canonical form of the action. Thus, the dressed coordinates can be also defined via a canonical transformation.
The calculus of the transition amplitudes has been performed using the dressed propagator, being obtained the basic formula which defines the *sum rules* presented in , then, the rules have been extended to describe other physical processes. In spite of the computation seems very difficult, the dressed coordinates allow to use the properties of the Hermite polynomials simplifying greatly the calculus.
On the other hand, it has been made an extensive use of the model given by Eq. (1) to study different physical situations, such as the quantum Brownian motion, decoherence and other related problems in quantum optics. In such context, it is interesting the computation of the reduced matrix density for the model (1) in the framework of dressed coordinates; the advances in such direction will be reported elsewhere.
### Acknowledgement
GFH (grant 02/09951-3) and RC (grant 01/12611-7) thank to FAPESP for full support. BMP thanks CNPq and FAPESP (grant 02/00222-9) for partial support.
## Appendix A The orthonormal matrix T=$`\left[t_\mu ^s\right]`$
Because the orthogonal character of the $`𝐓`$-matrix its components satisfy
$$\underset{\mu =0}{\overset{N}{}}t_\mu ^rt_\mu ^s=\delta _{rs},\underset{r=0}{\overset{N}{}}t_\mu ^rt_\nu ^r=\delta _{\mu \nu }$$
(67)
and from the Lagrangian (3) expressed in terms of the normal $`\left\{Q_r\right\}`$ we get other important relation
$$\underset{\mu =0}{\overset{N}{}}\overline{\omega }_\mu ^2t_\mu ^rt_\mu ^s\underset{k=1}{\overset{N}{}}\eta \omega _kt_0^rt_k^s\underset{k=1}{\overset{N}{}}\eta \omega _kt_0^st_k^r=\mathrm{\Omega }_r^2\delta _{rs}$$
(68)
where $`\overline{\omega }_0^2`$ have been defined in (4) and $`\overline{\omega }_k^2=\omega _k^2`$ . Using the equations above we can show the following sum
$$\overline{\omega }_0^2=\underset{s=0}{\overset{N}{}}\mathrm{\Omega }_s^2\left(t_0^s\right)^2,\eta \omega _k=\underset{s=0}{\overset{N}{}}\mathrm{\Omega }_s^2t_k^st_0^s,\omega _k^2=\underset{s=0}{\overset{N}{}}\mathrm{\Omega }_s^2\left(t_k^s\right)^2$$
(69)
and also compute the elements of the $`𝐓`$-matrix
$$t_k^s=\frac{\eta \omega _k}{\omega _k^2\mathrm{\Omega }_s^2}t_0^s,t_0^s=[1+\eta ^2\underset{k=1}{\overset{N}{}}\frac{\omega _k^2}{\left(\omega _k^2\mathrm{\Omega }_s^2\right)^2}]^{\frac{1}{2}}$$
(70)
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# Low Power Compact Radio Galaxies at High Angular Resolution
## 1 Introduction
Radio sources come in different well known and understood morphologies. Extended (kiloparsec scale) radio galaxies are divided between FR I and FR II radio galaxies, according to their morphology and radio power (Fanaroff & Riley 1974). Ledlow & Owen (1996) showed that the separation between these two classes is related to both the total radio power and the optical magnitude of the host galaxy. Ghisellini & Celotti (2001) showed how this separation can be interpreted in terms of accretion rates below or above a critical upper limit for maintaining an optically thin advection dominated accretion flow (ADAF, Narayan & Yi 1995) regime.
Compact (sub-kiloparsec) radio sources are also divided on the basis of different morphologies. Some of them appear as compact radio sources because of projection effects; this includes BL Lacertae objects (e.g., Mkn 501, Giroletti et al. 2004) and flat spectrum radio quasars (FSRQ), i.e., sources oriented at small angles with respect to the line of sight. Their compactness is due to projection effects, and to strong relativistic effects which give rise to one-sidedness, superluminal motions, and high brightness temperatures.
Other high power radio sources appear intrinsically small and are not affected by relativistic effects, since we do not see beamed relativistic jets but simply the regions of interaction between jets and the ISM (hot spots), which are advancing at velocities of the order of 0.2c (Owsianik & Conway 1998; Polatidis et al. 2002; Giroletti et al. 2003). The Compact Symmetric Objects class (CSOs, see e.g. Wilkinson et al. 1994; Gugliucci et al. 2005) is a good example of this population: it is composed of small sources ($`<1`$ kpc), with emission on both sides of the central engine. On the basis of kinematics as well as spectral arguments, these objects are interpreted as young radio galaxies with ages $`10^4`$ years and are expected to evolve into kiloparsec scale radio galaxies (Owsianik & Conway 1998; Murgia et al. 1999; Giroletti et al. 2003). The number of these sources appears, however, to be too high with respect to the general population of radio galaxies (Fanti et al. 1990; O’Dea & Baum 1997), therefore it has been suggested that their radio power decreases and/or expansion slows with time to account for the observed population of giant radio galaxies (Kaiser et al. 1997; Blundell et al. 1999; Fanti & Fanti 2003; Lara et al. 2004).
In catalogues of radio sources selected at low frequency there is another class of radio sources whose properties are not yet well known. We have named these low power compact (LPC) radio sources. Most of these sources do not have a flat radio spectra and show a moderately steep spectral index. Moreover, their host galaxies do not show signatures of strong nuclear activity in the optical and X-ray bands. X-ray luminosities, or upper limits, are typically of the order 10<sup>40</sup> erg s<sup>-1</sup> (e.g. Canosa et al. 1999). In the optical, central compact cores are seldom visible; and even when they are detected, these cores show a similar nature to the synchrotron optical cores of FR I radio galaxies rather than signatures of thermal emission from an efficiently radiating accretion disk (Capetti et al. 2002). The radio powers are typically below 10<sup>25</sup> W Hz<sup>-1</sup> at 1.4 GHz. Little variability is present in the radio fluxes, although this can be accounted for by the paucity of measurements.
The small size of LPCs could be directly related to the low radio power: the central AGN has insufficient power to drive the relativistic jet out of the dense ISM present in the central regions of the host galaxy. This could be the case for sources like NGC 4278 (Giroletti et al. 2005) and may constitute the link between radio loud and radio quiet AGN such as Seyfert galaxies. However, most of these low power compact radio sources have an intrinsic radio power in the same range as that of low power giant FR I radio galaxies. In this case, the reason for the compactness of these sources is unclear. It could be again any of the previous physical reasons: geometrical-relativistic effects, low age, or frustration by a denser than average ISM. It could be that all of these effects are present, and in addition some sources may even be prematurely dying (Marecki et al. 2003; Kunert-Bajraszewska et al. 2004; Gugliucci et al. 2005).
All these possibilities are interesting: in the case of geometric-relativistic effects, we would have a population of low power BL-Lacs where relativistic jets are present in spite of the low power of the AGN; in the case of young radio sources, we could test if their number is in agreement with current populations and discuss whether they will become giant radio galaxies or not. The LPC sources could be related to CSOs, but at lower radio power with respect to most CSOs studied up to now. Finally, frustration has been excluded as an explanation to the origin of CSOs, but could apply in some LPC sources, so it is important to measure kinematic ages for symmetric LPC sources.
To investigate these points we have selected from the Bologna Complete Sample (BCS, Giovannini et al. 2005) compact ($`<`$ 10 kpc) radio sources with a total low frequency radio power $`P_{408\mathrm{MHz}}<10^{25.5}`$ W Hz<sup>-1</sup> whose sub-parsec structure was not well defined. We are studying their properties with Very Long Baseline Array (VLBA) and high frequency, high resolution Very Large Array (VLA) observations. Two sources fulfilling the same requirements (the CSO 4C 31.04 and NGC 4278) have been already discussed in published papers (Giroletti et al. 2003, 2005). Here we present new results for 5 radio galaxies: 0222+36, 0258+35, 0648+27, 1037+30, and 1855+37 (see Table 1). The remaining compact sources will be observed in the near future and discussed in light of the statistical properties of the BCS.
The paper is laid out as follows: in §2 we give the details of our observations, in §3 we present the results about single sources, and in §4 we discuss the properties of our subsample with respect to other classes of compact sources; we present our main conclusions in §5. Throughout this paper, we make use of H<sub>0</sub> = 70 km sec<sup>-1</sup> Mpc<sup>-1</sup>, $`\mathrm{\Omega }_M`$ = 0.3 and $`\mathrm{\Omega }_\mathrm{\Lambda }`$ = 0.7. Spectral indices are defined such that $`S(\nu )\nu ^\alpha `$.
## 2 Observations and Data Reduction
### 2.1 VLA observations
VLA observations of five objects (0222+36, 0258+35, 0648+27, 1037+30, and 1855+37) were obtained in three observing runs in 2003 June 28 and 29 and 2003 July 10. The array was in “A” configuration (maximum baseline 35.4 km) and the observing frequencies were 8.4 and 22 GHz. Standard observing schedules for high frequency observations were prepared, including scans to determine the primary reference pointing, and using a short (3 s) integration time and fast switching mode (180 s on source, 60 s on calibrator) for K band (22 GHz) scans. Primary, amplitude, and phase calibrators for each run and for each source are given in Table 2. Post-correlation processing and imaging were performed with the NRAO Astronomical Image Processing System (AIPS). Parameters of different images are reported in Table 3.
### 2.2 VLBA observations
The 5 sources are also part of a sample of 15 objects that have been observed in phase reference mode with the VLBA, in order to study the parsec scale structure of faint radio galaxies. Observations were done in two separate runs on 2003 August 07 (BG136A, 24 hrs for 12 sources) and 2003 August 30 (BG136B, 6 hrs for 3 sources). We discuss here data for the 5 LPC sources only; they have all been observed in segment A of the experiment; the VLBA structure for the other 10 sources will be discussed in a future paper.
We observed in full polarization (RCP and LCP) with two IFs (central frequency 1659.49 MHz and 1667.49 MHz). We recorded 16 channels per frequency, for a total aggregate bit rate of 128 Mbs. Each pointing on a target source was bracketed by a calibrator scan in a 5 minutes duty cycle (3 minutes on source, 2 minutes on the calibrator). Two groups of (typically) 11 cycles were executed for each source at different hour angles, resulting in a total of about 66 minutes per target, with good coverage of the $`(u,v)`$plane. Calibrators were chosen from the VLBA calibrators list to be bright and close to the source; we report in Table 4 the list of the selected calibrator and its separation for each source discussed here. Short scans on strong sources (4C39.25, J0237+2848) were interspersed with the targets and the calibrators as fringe finder sources.
The correlation was performed in Socorro and the initial calibrations were done with AIPS. Scans on J0237+2848 were used to remove IF-dependent delays and phase offsets. After applying models for the electron content of the ionosphere as measured by the Jet Propulsion Laboratory (JPL) on the observing dates, we performed a global fringe fitting on all calibrators and applied the solutions to the targets with a two-point interpolation. $`RL`$ delay differences were determined and removed with the VLBACPOL procedure using 4C39.25. We then produced images of the calibrators in order to obtain and apply more accurate phase and gain corrections to the sources; we also determined from preliminary maps the absolute position of the sources, which we used thereafter. We give the coordinates of the core candidate of each source in Table 4.
The final calibrated single-source datasets were imported into Difmap for imaging and self-calibration. We produced images with natural and uniform weights, and summarize in Table 3 the significant parameters for the final images.
## 3 Results
All sources are detected in our VLA data at 8.4 GHz. The source structures are successfully resolved, allowing a study of different components. At 22 GHz, we have high signal to noise detections for 0222+36, 0258+35, and 0648+27; in the source 1037+30 we detect only the core and N-W hot spots; finally, 1855+37 is completely resolved and not even a faint nuclear source was detected. In sources 0222+36, 0258+35, and 0648+27 the lack of short baselines in 22 GHz images does not affect flux density measurements, therefore a spectral index comparison between 8.4 and 22 GHz is meaningful. We give in Table 5 the main parameters for different sources and subcomponents, such as flux densities and spectra.
The VLBA phase referenced images in total intensity are available for all sources except 1855+37, which is not detected. On the mas scale a resolved structure is visible in 0222+36, 0258+35, and 0648+27, while in 1037+30 we detect a hot spot in addition to the core, $`0.8\mathrm{}`$ north-west of it. Thanks to successful phase-referencing, our data yield accurate positional information which unambiguously identifies the nuclear component in these complex sources. Results are reported in Table 6, including spectral information confirming the core identification. The core spectral indices are computed between 1.6 and 22 GHz; for 0222+36 only, we consider the total flux density detected at 1.6 GHz by the VLBA, which extends over a region smaller than the beam of the VLA at 22 GHz. For other sources, secondary components are well separated in the VLA images (0258+35, 1037+30), or their flux density is negligible (0648+27). We also place upper limits on the fractional polarization. In fact, all images in polarized intensity are purely noise-like, with $`3\sigma `$ levels of approximately 0.25 mJy/beam. Since the sources are weak, these provide weak constraints of approximately $`P<10\%`$ at the peak position (except for 0222+36, see 3.1).
Finally, we also collected integrated flux density measurements from the NASA Extragalactic Database (NED), as well as from radio surveys including the Northern VLA Sky Survey (NVSS), the Westerbork Northern Sky Survey (WENSS), and the VLA Low-frequency Sky Survey (VLSS). These data typically cover the frequency range between 74 MHz and 5 GHz. Our new measurements at shorter wavelengths allow us to obtain integrated spectra spanning two and a half orders of magnitude in frequency. We show the integrated spectra for each source in the bottom right panels of Figs. 1, 4, 5, 6, 8, with best-fit continuous injection models (Murgia et al. 1999) overlaid. Detailed results for each source are presented in the following subsections.
### 3.1 0222+36
This source at $`z=0.0334`$ ($`1\mathrm{}=0.66`$ kpc) has a flux density of 337 mJy at 408 MHz, corresponding to a monochromatic total power of 10<sup>23.9</sup> W Hz<sup>-1</sup>. Previous VLA observations at 1.4 GHz reveal a slightly resolved morphology with a flat spectrum core surrounded by a halo extended over $`8\mathrm{}`$ (Fanti et al. 1986). On parsec scales, EVN and low resolution VLBA observations show an unresolved component with a flux density of $``$ 100 mJy at 5 GHz (Giovannini et al. 2001).
Our 8.4 GHz VLA observations resolve the structure of 0222+36 into a core and two components on either side (see Table 5). At 22 GHz the core is the dominant structure and the two lobes are faint with an S shaped structure (see Fig. 1).
The phase referenced VLBA image detects 102 mJy of flux density at 1.6 GHz in total intensity. No polarized flux is detected, which corresponds to a limit on the fractional polarization of $`<0.4\%`$ at the core position. The source is two sided, with jets emerging in opposite directions along the north-south axis. The unresolved central core has a flux density of 58.2 mJy. Both jets have a slightly bent path and become aligned with the kpc scale main axis at $`20`$ mas from the core. The resolved structure and the low core dominance are in agreement with the moderately steep spectrum of the VLA core ($`\alpha _{8.4}^{22}=0.69\pm 0.02`$; see Tab. 5).
Both the parsec and kiloparsec scale morphology strongly suggests that the source is oriented near the plane of the sky. At 10 mas from the peak, the jet/counter-jet brightness ratio is $`R=B_J/B_{CJ}=1.3`$ ($`B_J=12.0`$mJy and $`B_{CJ}=9.5`$mJy, with the main jet being the one pointing north); this corresponds to $`\beta \mathrm{cos}\theta `$ 0.05 which implies $`\theta >85^{}`$ if $`\beta >0.6`$. Therefore, the source is not affected by orientation effects and it has to be intrinsically small. If we consider the largest angular extent of the source and deproject it with an angle of $`85^{}`$ we derive an intrinsic size of approximately 5.4 kpc.
It is intriguing to try and understand the nature of the extended radio emission surrounding the central components in the form of a $`10`$ kpc halo. This halo is readily visible at low frequency (Fanti et al. 1986) but it is completely resolved in our images. For this reason, we reanalyzed VLA archive data at 1.4 GHz (D array) and 5 GHz (in A + B configuration). The 1.4 GHz D array data show that we are not missing a low brightness region more extended than the halo, and the 5 GHz data allow us to derive the spectral index of the halo region. We show in Fig. 2 the 5 GHz A+B array image.
From a comparison of images at different resolution (see Fig. 1 and 2), we do not see a clear connection between the inner structure (VLBA and high frequency VLA) and the extended halo; on the other hand, the VLBA structure is well connected to the source morphology as seen in the 22 and 8 GHz images. We consider it unlikely that the halo is due to the presence of extended lobes along the line of sight because this would imply a large bend between the small and large scales with considerable fine tuning required to produce such a symmetric, uniform, and circular extended emission.
We suggest that the halo could be due to the diffusion of relativistic particles and magnetic field around the source, similar to the extended halos seen in some spiral galaxies (Hummel et al. 1991). To better investigate this point, we consider the radio spectrum of different components. In Table 7 we report flux density measurements between 74 MHz and 22.5 GHz; the data are also shown in Fig. 3, along with a best-fit continuous injection model obtained using the Synage program (Murgia & Fanti 1996). The total spectrum (bottom right panel of Fig. 1) is complex but in good agreement with the sum of spectra of different components (Fig. 3). It is also worthwhile to remember that some variability could be present and that the data at our disposal were not taken simultaneously. However, we do not find any obvious inconsistency in the data, and based on this fact we argue that the variability is insignificant.
We further point out that the core emission refers to the sub-arcsecond core since we do not have multi-frequency VLBI data. For this reason, at 1.6 GHz we used the total correlated flux in our VLBA data for the core flux density. This value (102 mJy) is lower than the core flux density reported by Fanti et al. (1986), which however does definitely include a contribution from the lobes. The resulting spectrum for the sub-arcsec core shows a turnover with a maximum of 200 mJy at 3.3 GHz. Assuming that the turnover is produced by synchrotron self-absorption, we estimate, following Marscher (1987), the average magnetic field in the sub-arcsecond core to be in the range 2 – 9 $`\times `$ 10<sup>-2</sup> Gauss.
In our VLA images, both lobes show a steep spectrum between 8 and 22 GHz ($`\alpha =0.93\pm 0.04`$ and $`0.97\pm 0.04`$ in the NE and SW lobe, respectively). In order to obtain a better estimate of the physical parameters in the lobes, we also consider lower frequency data. These data have been obtained by subtracting the estimated flux density of the sub-arcsec core. For this reason we give in Fig. 3 and Table 7 the spectrum of the two lobes together. This explains the larger dispersion and uncertainty of this result. We estimate the average equipartition magnetic field in the two lobes with standard assumptions: we consider a frequency range from 10 MHz to 100 GHz, an equal amount of energy in heavy particles and in electrons ($`k=1`$), and that relativistic particles and magnetic fields occupy the same volume ($`\varphi =1`$). Under these assumptions, we find B<sub>eq</sub> = 1.3 $`\times `$ 10 <sup>-4</sup> Gauss in the lobes and a break frequency $`\nu _{\mathrm{br}}9`$ GHz.
On larger scales, we derive a flux density of the halo from the archival data in A and B configuration at 5 GHz (see the image obtained combining A+B data in Fig. 2). Similarly to the case of the lobes, we derive the equipartition magnetic field in the extended halo. The spectrum of this region is straight and steep ($`\alpha `$ 1.2), with a low break frequency $`\nu _{\mathrm{br}}0.3`$ GHz. In this region we have B$`{}_{\mathrm{eq}}{}^{}6.9\times 10^6`$ Gauss.
From the spectral index information and the equipartition magnetic field, we estimate the radiative age in the lobes and in the halo region. We find that the lobes are 4.5 $`\times `$ 10<sup>5</sup> yrs old and the extended halo is $`1.3\times 10^8`$ yrs old. Therefore, we have in this source a young structure (inner lobes) surrounded by an older symmetric and diffuse region. Since there is evidence of a change of the jet direction in the inner young region, we can speculate that the jet orientation is rotating because of instabilities of the AGN. This could explain the small size of the source because an unstable jet did not allow the growth of a large scale radio galaxy. In this case the old round halo could be due to the diffusion of radio emission during the orbit of the inner structure. From VLBA and high resolution VLA images we can estimate that the jet rotated by $``$ 40 in 5 $`\times `$ 10<sup>5</sup> yrs so that in 10<sup>8</sup> yrs (the halo age) it could have done many complete orbits. Higher sensitivity A+B+C data at 8.4 GHz might be able to follow the radio structure from the lobes to the halo to test this hypothesis.
We can also compare our estimate of the spectral age of the lobes to a reasonable dynamic estimate based on their size. Assuming that the two lobes separation velocity is $``$ 0.1 c, we derive a dynamic age $`2\times 10^4`$ yrs. This age is 10 times lower than the radiative age, but such differences have been found in other CSOs (e.g. 4C31.04, Giroletti et al. 2003) and could be a result of the many assumptions going into both the spectral and kinematic age estimates (e.g., constant advance velocity and equipartition conditions).
### 3.2 0258+35 (NGC 1167)
At $`z=0.0165`$ ($`1\mathrm{}=0.34`$ kpc), this source has a radio power at 408 MHz of $`10^{24.37}`$ W Hz<sup>-1</sup>. It is optically classified as a Seyfert 2 galaxy (Ho et al. 1997), and in the radio it has been previously studied and classified as a compact steep spectrum (CSS) source by Sanghera et al. (1995), even though the total spectrum is only moderately steep: $`\alpha _{0.08}^{22}=0.54`$. In contrast with most CSS sources, 0258+35 shows plume-like lobes without prominent hot-spots. Previous studies include the global VLBI image at 5 GHz in Giovannini et al. (2001), which reveals a compact component with a short jet-like structure and an excess of emission on the short spacings.
Our VLA images are shown in Fig. 4 (top panels). The 8.4 GHz image shows a morphology in good agreement with previous images. Thanks to the high sensitivity of the VLA at this frequency, an extended diffuse emission surrounding the inner radio structure is detected. In the 22 GHz image, thanks to the better angular resolution, the source peak is resolved into two components, one of which is fainter and more central, while the other is stronger and located at the beginning of the SE lobe.
Our phase-referenced VLBI data show a faint compact component identified with the nuclear source and an extended “blob” at $`0.1\mathrm{}`$. As shown by the cross overlaid on the VLA images, the compact parsec scale core is located in the central component at the base of the south-east jet-like feature. This component is easily identifiable in the 22 GHz image, while in the 8 GHz map the core is confused with a bright jet component. Besides the faint compact core ($`S_\mathrm{c}=7.6`$ mJy), the largest fraction of flux density (240 mJy) in the VLBI image at 1.6 GHz is contained in a bright diffuse region coincident with the peak of the VLA images. The flux density and size of the components are given in Table 6. We wonder about the nature of this region because of its morphology and the lack of a connection with the parsec-scale core. It could be the result of a burst of activity of the central AGN, in which case some moderate Doppler boosting is necessary to explain the absence of a similar counterjet on the opposite side. From the jet to counter-jet ratio in the 22 GHz VLA map measured at $`0.5\mathrm{}`$ from the core, we find that the source has to be oriented at an angle $`<`$ 70; if the source has subarcsecond relativistic jets, the orientation is in the range $`60^{}\theta 70^{}`$, in agreement with the optical classification of the parent galaxy as a Seyfert 2. From the VLBI image and using the knot brightness we find $`R=B_J/B_{CJ}7.5/0.1358`$, resulting in $`\beta \mathrm{cos}\theta >0.67`$, and a viewing angle $`40^{}\theta 50^{}`$, if $`\beta 0.9`$. However, we cannot rule out the possibility that the knot was produced in a peculiar episode and the jet/counterjet ratio is not meaningful.
The 22 GHz VLA image has about the same resolution ($`80`$ mas) as the combined EVN+MERLIN 1.6 GHz map shown in Sanghera et al. (1995), and we can estimate the spectral index for the core ($`\alpha 0.0`$) and main jet component ($`\alpha 0.4`$). A flat core spectral index is confirmed by our VLA images, albeit a slight misalignment between the two images prevents an unambiguous core identification only from our data. The spectrum is still flat in the central region and gradually steepens to become almost constant (0.6) in the inner bright jet-like structure in both lobes. The surrounding diffuse emission is steep: 1.0 – 1.5. The integrated spectral index at high frequency is $`\alpha _{8.4}^{22}=0.88\pm 0.01`$.
Since the knot is well resolved in our VLBA image we derive its opening angle. We measure an angle of 26, which is close to the intrinsic opening angle if the source is oriented near the plane of the sky. A free expanding jet is expected to show an intrinsic opening angle of $`1/\gamma `$ (see, e.g. Salvati et al. 1998); under this assumption, we estimate a Lorentz factor $`\gamma `$ 2.4 which corresponds to $`\beta `$ 0.9.
The source structure suggests that in the data at 5 GHz (Giovannini et al. 2001) the main peak is produced by the nuclear source with the beginning of the main jet also visible. In this case the VLBI core has an inverted spectrum, free-free or self-absorbed (26 mJy at 5 GHz and 7.5 mJy at 1.6 GHz).
We have estimated the average equipartition magnetic fields in the source, assuming an uniform brightness in the source volume. This can be considered a good approximation, since the compact structures (core and knot) are only $`20\%`$ of the total flux density at 5 GHz. We estimate B<sub>eq</sub> $``$ 9 $`\times `$ 10<sup>-5</sup> Gauss. With this estimate and using the break frequency found in the total spectral index distribution (4.6 GHz) we estimate an age of 9 $`\times `$ 10<sup>5</sup> yrs for this source. Of course, this is an average age and we expect that external diffuse regions are older with respect to the innermost region.
In light of these results, we speculate that this source might not grow to become a kiloparsec scale radio galaxy. No final hot spots demarcating the ends of the jets are visible and the source structure appears to strongly interact with the ISM as shown by the large bending of the arcsecond structure of the SE lobe and the presence of a surrounding low brightness extended structure in the VLA images. We note that in this source the inner jet direction appears to be constant, moreover no amount of bending is visible in the NW lobe, therefore the source structure on the large scale should be related not to the inner BH motion but to interaction with the ISM. We note that the estimated radiative age should imply a larger source size even allowing for a low lobe advance velocity.
### 3.3 0648+27
This object ($`z=0.0414`$, corresponding to 0.82 kpc/″, $`P_{408\mathrm{MHz}}=10^{24.02}`$ W Hz<sup>-1</sup>) is only slightly extended at the lowest frequencies (Parma et al. 1986). It was resolved into a double source extended about 1″ with VLA observations at 8.4 GHz by Morganti et al. (2003); they also detect a large amount of Hi using the Westerbork Synthesis Radio Telescope (WSRT, $`M_{\text{H}\text{i}}=1.1\times 10^{10}M_{}`$). However, Morganti et al. (2003) are not able to identify a core and consequently interpret the radio structure in terms of a pair of symmetric lobes.
Our 8.4 GHz VLA data (Fig. 5, top left panel) confirm the structure as a double, with flux densities of 19 and 14 mJy in the northern and southern components, respectively, in agreement with the data from Morganti et al. (2003). However, both components are resolved at 22 GHz (see top right panel in Fig. 5), and a compact feature emerges in the north with a flux density of 2.5 mJy. Its small size ($`<0.07\mathrm{}`$) and flatter spectral index suggest that this component is actually the core, with emission on either side. See Table 5 for a list of our data and spectral index measurements for the two lobes.
Our phase referenced VLBI data lend strong support to this scenario. In fact, the emission in the VLBA image is located in the vicinity of the VLA peak in the northern lobe. The total flux density is only 12.8 mJy, with a peak of 4.4 mJy/beam. A faint jet-like structure is visible to the south-east, although the signal-to-noise is very poor and it could be spurious. In any case the difference between the total correlated flux and the peak flux density in the VLA image suggests the presence of some extended emission on intermediate scales. We note that the spectral index between the 22 GHz VLA unresolved component and the total VLBA correlated flux is $`\alpha _{1.6}^{22}=0.47\pm 0.03`$ (see Table 6). It is possible that this component is a jet knot and the core is at the extreme northern end. The core might be free-free or self-absorbed at 1.6 GHz, so that the VLBA does not pick out the center of activity.
The source flux density is dominated by the extended emission. The total spectrum of the source between 325 MHz and 22 GHz has an index $`\alpha 0.8`$, with a hint of steepening at high frequency ($`\alpha _{8.4}^{22}=1.22\pm 0.02`$). A measurement at 5 GHz ($`S_{5\mathrm{GHz}}=213`$ mJy, Antonucci 1985) deviates significantly from the other measurements and was not included in the spectrum. If it is not due to an error, then it is difficult to understand this value.
The two lobes do not show any evidence for jet-like structure or the presence of hot spots. The lobe structure is relaxed with a relatively steep spectrum. The South lobe is at a greater distance with respect to the core, but we ascribe this asymmetry to a difference in the ISM and not to a relativistic effect. We note also that the lack of prominent jet structures in the parsec scale image suggests that no relativistic jet structure is present or that the source is on the plane of the sky. The faintness of the nuclear emission also argues against a Doppler boosted jet.
We have estimated the average equipartition magnetic field in this source to be H$`{}_{\mathrm{eq}}{}^{}95\times 10^6`$ Gauss. Using the break frequency (2.9 GHz) estimated from total flux density measures we can estimate a minimum age for this source of about 1 Myr ($`9.9\times 10^5`$ years). We expect that the external lobe regions are much older, confirming that this source is confined, and is expected to remain compact similar to NGC 4278 (Giroletti et al. 2005) despite its relatively higher total radio power and larger size. The connection between the small size of the radio emission and the presence of a major merger in this galaxy about 10<sup>9</sup> years ago, with the presence of a large amount of Hi in this galaxy (Morganti et al. 2003), is remarkable.
### 3.4 1037+30
Located at $`z=0.0911`$ ($`1\mathrm{}=1.70`$ kpc), 1037+30 has $`P_{408\mathrm{MHz}}=10^{25.37}`$ W Hz<sup>-1</sup>. It is only slightly resolved at 1.4 GHz (Fanti et al. 1986) and not detected in any VLBI observations to date (Giovannini et al. 2005). In the optical, 1037+30 is identified with the brightest galaxy in the cluster Abell 923.
Our VLA images are shown in the top panels of Fig. 6. The 8.4 GHz image reveals an edge-brightened structure, with complex sub-structures: jets, and lobes with hot spots. In the 22 GHz image only a point-like component, probably the core, and the resolved NW hot spot are evident. As in other objects, the accurate phase-referenced VLBI image provides confirmation for the identification of the core which appears as a faint central component in the 8.4 GHz image.
The south-eastern jet is clearly visible in the 8.4 GHz VLA image, and is both longer and better defined than the jet in the opposite direction. While the jet brightness and length suggest that the SE is the approaching side, the NW hot spot is much brighter (with a peak of 13 vs 2.3 mJy/beam at 8.4 GHz). It is possible that the ISM density is irregular and that the brightness asymmetry is due to interactions with an inhomogeneous ISM.
The NW hot spot has a peculiar morphology, with a very sharp edge. We interpret this structure in terms of a back-flow due to a strong interaction of the NW jet with the ISM. This could also explain the compressed appearance of the NW hot spot, its higher brightness, and the “bridge” connecting the NW hot spot with the SE hot spot. Since the jet to counterjet brightness ratio in the 8.4 GHz image is not high, we expect that the source is oriented at a large angle with respect to the line of sight. This orientation is in agreement with the non-detection of the jet at 22 GHz and at 1.6 GHz by the VLBA. This orientation is also consistent with the low core dominance, since relativistic jets at a large angles will have a Doppler factor $`<<`$ 1.
The VLBA data show a clear detection of a 4 mJy core (see Table 6). Faint, diffuse emission is detected on the shortest baselines, although it is not possible to image it properly. At the very limit of the VLBA field-of-view, we detect radio emission from the north west hot spot region. Our image (Fig.7) shows a resolved structure $`60`$ mas in size in agreement with the 22 GHz image. The total flux density at 1.6 GHz in this region is $`40`$ mJy. Deeper images are necessary to see the connection to the jet.
The size and morphology of this source are in agreement with the definition of CSO sources. The jet interaction with the ISM is well defined. We estimated a dynamical age assuming a lobe expansion velocity of 0.2c and find an age of $`4.5\times 10^4`$ yrs assuming for the size the distance from the core of the SE hot spot. We do not estimate for this source a synchrotron age: the complex radio structure makes any estimate of an average H<sub>eq</sub> unrealistic, and the contribution of different components to the total spectrum is difficult to weigh. The overall break frequency is however around 2.7 GHz.
We note that because of its total radio power we expect that 1037+30 should evolve into an extended FR I radio galaxy. This source, like 0116+31 (4C 31.04, Giroletti et al. 2003), is also one of a number of low power radio sources with a clear CSO morphology.
### 3.5 1855+37
Located at a redshift $`z=0.0552`$ ($`1\mathrm{}=1.07`$ kpc), 1855+37 has a total power $`P_{408\mathrm{MHz}}=10^{24.65}`$ W Hz<sup>-1</sup>. Compact in the NVSS, this object is resolved into a triple source with higher resolution VLA observations at 1.4 GHz (Fanti et al. 1986) and 5 GHz (Morganti et al. 1997), where it also shows polarization. We identify 1855+37 with the brightest member of the galaxy cluster CIZA J1857.6+3800, discovered on the basis of X-ray data by Ebeling et al. (2002) in the zone of avoidance. A discussion of the X-ray properties of 1855+37 on the basis of ROSAT PSPC data is presented in Canosa et al. (1999) and Worrall & Birkinshaw (2000).
Our observations at 8.4 GHz (Fig. 8) resolve the flux density of $`40`$ mJy into a weak, diffuse emission extended over $`7\mathrm{}`$. The structure is two sided, with a weak central component ($`S_{\mathrm{core}}=0.57`$ mJy) and large, faint lobes. The source is not detected at 22 GHz; no significant peak is found in this region in the VLBA image.
Given the low power of the core, the lack of visible jets, and the rather steep spectral index, it is possible that the activity in this source is fading away. Note that in all previous observations, typically at lower frequency (Fanti et al. 1986; Morganti et al. 1997), the source was largely dominated by the extended lobes. For this reason, it was not possible to pinpoint the location of the nuclear activity, whose extreme weakness has therefore gone unnoticed. Our 8.4 GHz observations reveal for the first time the exact position and flux density of the core, which is remarkably weak ($`S_{\mathrm{core}}/S_{\mathrm{tot}}=0.017`$). This also explains the lack of detection of nuclear activity in 22 GHz VLA data and in our VLBI observations. We note also that according to the correlation between the core and total radio power (see e.g. Giovannini et al. 2001) the nuclear power in this source is too low even assuming a large jet velocity and a huge Doppler de-boosting. We therefore propose this galaxy hosts a dying radio source. If the nucleus of this galaxy is ceasing its activity we would not expect that this source would significantly grow any larger than it is now ($`7`$ kpc).
We note that the present radio structure is very similar to that found in NGC 4874 (the brightest galaxy in the Coma cluster) by Feretti & Giovannini (1985). However, in NGC 4874 the presence of a core and of bright regions in the two lobes suggests ongoing nuclear activity.
The few data points available in the literature, combined with our new measurements at 8.4 GHz and limit at 22.5 GHz, yield a break frequency $`\nu _{\mathrm{br}}2.5`$ GHz. The equipartition magnetic field is $``$ 2.1 $`\times `$ 10<sup>-5</sup> Gauss and the corresponding radiative source age is 9.7 Myrs. Since this source is inside an X-ray cluster we expect that its small size with respect to the radiative age can also be due to confinement from the IGM gas (see also Worrall & Birkinshaw 2000). Also the source shape, similar to wide angle tail radio sources, could be due to relative motions between thermal gas and the source.
## 4 Discussion
We present new VLA and VLBA results for 5 Low Power Compact Radio Sources. Our high resolution images resolve all sources. The combined information of flat/inverted spectral indices ($`0.2\alpha _{1.6}^{22}0.5`$) and compactness unambiguously identify radio cores in four of them. In one case (1855+37), we have a compact component at 8.4 GHz, which is a reasonable core candidate; however, a detection at another frequency would be desirable to confirm the classification. Most sources are symmetric, at least on arcsecond scales; 0222+36 is also clearly symmetric on milliarcsecond scales. In general, the objects do not show any evidence for strong beaming effects (e.g. strong, one-sided jets, dominant cores). Thus, these sources are most likely close to the plane of the sky. We determine accurate values for their (projected) linear size, which must be, therefore, within a factor of a few from their intrinsic deprojected dimensions.
We show in Fig. 9 a radio power vs linear size diagram for the 95 sources in the Bologna Complete Sample (BCS). Radio power data are taken from the B2 survey at 408 MHz via the on-line VizieR service<sup>1</sup><sup>1</sup>1http://vizier.u-strasbg.fr/viz-bin/VizieR (Ochsenbein et al. 2000) and linear sizes from several works in the literature<sup>2</sup><sup>2</sup>2The largest part of the data is drawn from Fanti et al. (1987) and from “An Atlas of DRAGNs” for sources that are part of the 3CRR sample (see http://www.jb.man.ac.uk/atlas/); a few individual sources are better studied in other references: 0116+31 (4C 31.04, Giroletti et al. 2003), 0836+29B (Fanti et al. 1986), 1144+35 (Giovannini et al. 1999), 1217+29 (NGC 4278, Giroletti et al. 2005), 1257+28 (NGC 4874, Feretti & Giovannini 1985), 1557+26 (from the FIRST survey, Becker et al. 1995), 1652+39A (Mkn 501, Cassaro et al. 1999). The solid points denote the five LPC sources in the subsample considered in the present work with accurate linear size measurements. We confirm that this class of radio sources is characterized by a linear size $`<`$ 10 kpc. The starred points indicate the two smallest sources in the BCS ($`LS<0.1`$ kpc), one of which is a genuinely young CSO (4C 31.04, Giroletti et al. 2003), while the other is associated with a less active LINER galaxy (NGC 4278, Giroletti et al. 2005). From the distribution shown in the diagram of Fig. 9, we note that the small size of LPCs is not related to the low power of the radio source. In fact, there is also a large number of extended ($`LS>10^2`$ kpc) radio galaxies with a similarly low radio power.
Our new LPC sources are more extended than classical young radio sources like the CSOs (Gugliucci et al. 2005), and could represent the intermediate age sources in between CSOs and classic kpc scale radio galaxies. The youth scenario has the problem of the lack of objects filling in the gap between sources with age $`<10^4`$ yrs and $`>10^7`$ yrs; in particular, there seem to be too many bright young sources. However, this problem can be solved if the radio sources dim in luminosity as they grow up; for example, assuming equipartition conditions and expansion losses, Begelman (1996) reconciles the observational data with a model that predicts that the radio luminosity should decrease with size, roughly as $`P_RLS^{1/2}`$. According to the evolution scenario young radio sources are then expected to be small and quite luminous.
In order to better understand radio source evolution and the radio power vs linear size diagram, we need to know the properties of compact radio sources. It is also important to point out that the constraints posed by the number counts per linear size range may be misleading. In fact, a number of explanations have been invoked, for example contamination from core-jet Doppler-boosted sources (Tinti et al. 2005) and a significant fraction of either frustrated or short lived sources (Readhead et al. 1994; Gugliucci et al. 2005).
To investigate these possibilities we are studying in detail compact (LS $`<10`$ kpc) sources in our BCS, which are typically low power ones. If we consider the data presented here for 5 sources and include also the two other sources previously studied (4C 31.04, NGC 4278; Giroletti et al. 2003, 2005), we are presented with a variety of behaviors. As a matter of fact, the individual study of each source suggests that all cases are possible and we have evidence in a number of cases that sources are not able to grow because of an underpowered core or a jet instability, or because they are dying. On the other hand, we still have only a few cases in which there is clear evidence of a source which is growing, with significant interaction between the jet and the ISM. In particular, among sources studied here, we have:
– One source with evidence of a growing jet structure interacting with the ISM (1037+30) with an estimated kinematic age of 5 $`\times `$ 10<sup>4</sup> yrs.
– A peculiar source in which an old structure and a younger one coexist (0222+36, in which the estimated radiative age of the halo is $`10^8`$ yrs, while the lobes are $`10^410^5`$ yrs old). We interpret this structure as due to a jet instability which does not allow it to create large scale lobes. The old halo is permeated by relativistic electrons which slowly escape from the younger region (lobes) and diffuse into the surrounding volume.
– Three sources that are small because they are short lived or frustrated. In 0258+35, which has no hot spots and resembles a FR I radio galaxy, the VLA structure does not show any evidence of interaction between expanding jets and the ISM. The VLBA structure (a faint core and an isolated more powerful knot) suggests variable levels of activity. In 0648+27 and 1855+37 we find a low power nuclear region and little or weak extended emission, suggesting that they are possibly frustrated and/or dying. In the Hi rich 0648+27, we detected a nuclear source with the VLBA and the estimated radiative age is $``$ 10<sup>5</sup> yrs, while in 1855+37 the nuclear source is no longer active and the source radiative age is 10<sup>7</sup> yrs. In general, the dynamic ages estimated from their large-scale structure and assuming typical lobe advance velocities are about one order of magnitude lower than the spectral estimates. Unless there is some misleading assumption (e.g., equipartition conditions do not apply), the advance velocities ($`<<0.1c`$) have to be significantly lower than those measured in CSOs.
## 5 Conclusions
Despite the poor statistics, our preliminary studies indicate that multiple causes can produce sources in the LPC class. In addition to flat or inverted spectrum sources dominated by projection effects as BL-Lacs, a small size can stem from: youth (4C 31.04 and 1037+30), instabilities in the jets (in space, as in 0222+36, or time, as in 0258+35), frustration (0648+27), a premature end of nuclear activity (1855+37), or just a very low power core (NGC 4278). A more detailed discussion will appear in a future paper about the BCS. The study of the BCS, as a well defined complete sample, will allow us to make a statistical study to derive the source evolution. From the number of sources in different evolutionary stages it should be possible to estimate the duration and probability of different stages in radio source life.
###### Acknowledgements.
We thank the referee Dr. C. Stanghellini for a prompt and useful report. The National Radio Astronomy Observatory is a facility of the National Science Foundation operated under cooperative agreement by Associated Universities, Inc. This research has made use of the NASA/IPAC Extragalactic Database (NED) which is operated by the Jet Propulsion Laboratory, Caltech, under contract with NASA and of NASA’s Astrophysics Data System (ADS) Bibliographic Services. This material is based upon work supported by the Italian Ministry for University and Research (MIUR) under grant COFIN 2003-02-7534.
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# Statistical model selection methods applied to biological networks
## 1 Introduction
Network structures which connect interacting particles such as proteins have long been recognised to be linked to the underlying dynamic or evolutionary processesdor\_book ; albert\_RMP2002 . In particular the technological advances seen in molecular biology and genetics increasingly provide us with vast amounts of data about genomic, proteomic and metabolomic network structures ito\_PNAS2000 ; wagner\_MBE2001 ; qin\_PNAS2003 . Understanding the way in which the different constituents of such networks, — genes and their protein products in the case of genome regulatory networks, enzymes and metabolites in the case of metabolic networks (MN), and proteins in the case of protein interaction networks (PIN) — interact can yield important insights into basic biological mechanisms maslov\_S2002 ; yook\_proteomics2004 ; Agrafioti2005 . For example the extent of phenotypic plasticity allowed for by a network, or levels of similarity between PINs in different organisms presumably depend on topological (in a loose sense of the word) properties of networks.
Our analysis here focuses on the degree distribution of a network, i.e. the probability of a node to have $`k`$ connections to other nodes in the network. While it is well known that this does not offer an exhaustive description of network data, it has nevertheless remained an important characteristic/summary statistic of network data. Here we use $`\text{Pr}(k)`$ to denote a theoretical model for the degree distribution, or $`\text{Pr}(k;\theta )`$ if the model depends on an (unknown) parameter $`\theta `$ (potentially vector-valued), and $`\widehat{\text{Pr}}(k)`$ to denote the empirical degree distribution.
Many studies of biological network data have suggested that the underlying networks show scale-free behaviour bollobas\_handbook and that their degree distributions follow a power-law, i.e.
$$\text{Pr}(k;\gamma )=k^\gamma /\zeta (\gamma )$$
(1)
where $`\zeta (x)`$ is Riemann’s zeta-functions which is defined for $`x>1`$ and diverges as $`x1`$; for finite networks, however, it is not necessary that the value of $`\gamma `$ is restricted to values greater than 1.
These powerlaws are in marked contrast to the degree distribution of the Erdös-Rényi random graphs bollobas\_book1 which is Poisson, $`\text{Pr}(k;\lambda )=e^\lambda \lambda ^k/k!`$. The study of random graphs is a rich field of research and many important properties can be evaluated analytically. Such Poisson random networks (PRN) are characterized by most nodes having comparable degree; the vast majority of nodes will have a connectivity close to the average connectivity.
The term ”scale-free” means that the ratio $`\text{Pr}(\alpha k)/\text{Pr}(k)`$ depends on $`\alpha `$ alone but not on the connectivity $`k`$. The attraction of scale-free models stem from the fact that some simple and intuitive statistical models of network evolution cause powerlaw degree distribution. Scale-free networks are not, however, the only type of network that produces fat-tailed degree distributions.
Here we will be concerned with developing a statistically sound approach for inferring the functional form for the degree distribution of a real network. We will show that relatively basic statistical concepts, like maximum likelihood estimation and model selection can be straightforwardly applied to PIN and MN data. In particular we will demonstrate how we can determine which probability models best describe the degree distribution of a network. We then apply this approach in the analysis of real PIN data from five model organisms and MN data. In each case we can show that the explanatory power of a standard scale-free network is vastly inferior compared to models that take the finite size of the system into account.
## 2 Statistical tools for the analysis of network data
Here we are only concerned with methods aimed at studying the degree distribution of a network. In particular we want to quantify the extent to which a given functional form can describe the degree distribution of a real network. Given a probability model (e.g. power-law distribution or Poisson distribution) we want to determine the parameters which describe the degree distribution best; after that we want to be able to distinguish which model from a set of trial model provides the best description. Here we briefly introduce the basic statistical concepts employed later. These can be found in much greater detail in most modern statistics texts such as davison\_book1 . Tools for the analysis of other aspects of network data, e.g. cluster coefficients, path length or spectral properties of the adjacency matrix will also need to be developed in order to understand topological and functional properties of networks.
There is a well established statistical literature that allows us to assess to what extent data (e.g. the degree distribution of a network) is described by a specific probability model (e.g. Poisson, exponential or powerlaw distributions). Thus far, determining the best model appears to have been done largely by eye dor\_book and it is interesting to apply a more rigorous approach, although in some published cases maximum likelihood estimates were used to determine the value of $`\gamma `$ for the scale-free distribution.
### 2.1 Maximum likelihood inference
Since we only specify the marginal probability distribution, i.e. the degree distribution, we take a composite likelihood approach to inference, and treat the degrees of nodes as independent observations. This is only correct in the limit of an infinite sized network and finite sized sample ($`n<<N`$, where $`N`$ denotes network size and $`n`$ the sample size). Composite likelihood methods are becoming increasingly popular in cases for which the full likelihood is difficult to specify and/or the full likelihood is intractable to calculate numerically. In our case the full likelihood is difficult to specify. Reference cox\_reid provides an overview of composite likelihood methods.
For a given functional form or model $`\text{Pr}(k;\theta )`$ of the degree distribution we can use maximum likelihood estimation applied to the composite likelihood in order to estimate the parameter which best characterizes the distribution of the data. The composite likelihood of the model given the observed data $`K=\{k_1,k_2,\mathrm{},k_n\}`$ is defined by
$$L(\theta )=\underset{i=1}{\overset{n}{}}\text{Pr}(k_i;\theta ),$$
(2)
and taking the logarithm yields the log-likelihood
$$\text{lk}(M)=\text{lk}(\theta )=\underset{i=1}{\overset{n}{}}\mathrm{log}(\text{Pr}(k_i;\theta )).$$
(3)
The maximum likelihood estimate (MLE), $`\widehat{\theta }`$, of $`\theta `$ is the value of $`\theta `$ for which Eqns. (2) and (3) become maximal. For this value the observed data is more probable to occur than for any other parameters.
Here the maximum likelihood framework is applied to the whole of the data. This means that in fitting a curve —such as a powerlaw $`k^{\widehat{\gamma }}/\zeta (\widehat{\gamma })`$, where $`\widehat{\gamma }`$ denotes the MLE of the exponent $`\gamma `$— data for all $`k`$ is considered. If a powerlaw-dependence where to exist only over a limited range of connectivities then the global MLE curve may differ from such a localized power-law (or equivalently any other distribution).
### 2.2 Model selection and Akaike weights
We are interested in determining which model describes the data best. For non-nested models (as are considered here, e.g. scale-free versus Poisson) we cannot use the standard likelihood ratio test but have to employ a different information criterion to distinguish between models: here we use the Akaike-information criterion (AIC) to choose between different models akaike\_proc1983 ; burnham\_book . The AIC for a model $`\text{Pr}(k;\theta )`$ is defined by
$$\text{AIC}=2(\text{lk}(\widehat{\theta })+d),$$
(4)
where $`\widehat{\theta }`$ is the maximum liklihood estimate of $`\theta `$ and $`d`$ is the number of parameters required to define the model, i.e. the dimension of $`\theta `$. Note that the model is penalized by $`d`$. The model with the minimum AIC is chosen as the best model and the AIC therefore formally biases against overly complicated models. A more complicated model is only accepted as better if it contains more information about the data than a simpler model. (It is possible to formally derive the AIC from Kohn-Sham information theory.) Other information criteria exist, e.g. the Bayesian information criterion (BIC) offers a further method for penalizing more complex models (i.e. those with more parameters) unless they have significantly higher explanatory power (see burnham\_book for details about the AIC and model selection in statistical inference). In order to compare different models we define the relative differences
$$\mathrm{\Delta }_j^{\text{AIC}}=\text{AIC}_j\underset{j}{\mathrm{min}}(\text{AIC}),$$
(5)
where $`j`$ refers to the $`j`$th model, $`j=1,2,\mathrm{},J`$, and $`\mathrm{min}_j`$ is minimum over all $`j`$.This in turn allows us to calculate the relative likelihoods (adjusted for the dimension of $`\theta `$) of the different models, given by
$$\mathrm{exp}(\mathrm{\Delta }_j^{\text{AIC}}/2).$$
(6)
Normalizing these relative likelihoods yields the so-called Akaike weights $`w_j`$,
$$w_j=\frac{\mathrm{exp}(\mathrm{\Delta }_j^{\text{AIC}}/2)}{_{j=1}^J\mathrm{exp}(\mathrm{\Delta }_j^{\text{AIC}}/2)}.$$
(7)
The Akaike weight $`w_j`$ can be interpreted as the probability that model $`j`$ (out of the $`J`$ alternative models) is the best model given the observed data and the range of models to choose from. The relative support for one model over another is thus given by the ratio of their respective Akaike weights. If a new model is added to the $`J`$ existing models then the analysis has to be repeated. The Akaike weight formalism is very flexible and has been applied in a range of context including the assessment of confidence in phylogenetic inference strimmer\_proysoc2002 . In the next section we will apply this formalism to PIN data from five species and estimate the level of support for each of the models in table 1.
### 2.3 Goodness-of-fit
In addition to the AIC or similar information criteria we can also assess a model’s performance at describing the degree distribution using a range of other statistical measures. The Kolmogorov-Smirnoff (KS)davison\_book1 and Anderson-Darling (AD) Anderson1952 ; Anderson1954 goodness-of-fit statistics allow us to quantify the extent to which a theoretical or estimated model of the degree distribution describes the observed data. The former is a common and easily implemented statistic, but the latter puts more weight on the tails of distributions and also allows for a secular dependence of the variance of the observed data on the argument (here the connectivity $`k`$). They KS statistic is defined as
$$D=\mathrm{max}|\widehat{P}(k)P(k)|,$$
(8)
where $`\widehat{P}(k)`$ and $`P(k)`$ are the empirical and theoretical cumulative distribution functions, respectively, for a node’s degree i.e. $`P(k)=_{i=1}^k\text{Pr}(i)`$ and $`\widehat{P}(k)=_{i=1}^k\widehat{\text{Pr}}(i)`$. If $`P(k)`$ depends on $`\theta `$, $`P(k)`$ is substituted by $`P(k;\widehat{\theta })=_{i=1}^k\text{Pr}(i;\widehat{\theta })`$, the estimated cumulative distribution. This statistic is most sensitive to differences between the theoretical (or estimated) and observed distributions around the median of the data, i.e. the point where $`P(k)0.5`$. Given that we will also be considering a number of fat-tailed distributions this is somewhat unsatisfactory and we will therefore also use the AD statistic (Anderson and Darling discussed a number of statistics Anderson1952 ; Anderson1954 ) which is defined as
$$D^{}=\mathrm{max}\frac{|\widehat{P}(k)P(k)|}{\sqrt{P(k)(1P(k))}}$$
(9)
(again $`P(k)`$ might be substituted for $`P(k;\widehat{\theta })`$).
We can use these statistics for two purposes: first, we can use them to compare different trial distributions as the ”best” distribution should have the smallest value of $`D`$ and $`D^{}`$, respectively. Second, we can use these statistics to determine if the empirical degree distribution is consistent with a given theoretical (or estimated) distribution.
To evaluate the fit of a model, $`p`$-values can be calculated for the observed values of $`D`$ and $`D^{}`$ using a parametric boot-strap procedure using the estimated degree distribution: for a network with $`N`$ nodes we repeatedly sample $`N`$ values at random from the maximum likelihood model, $`\text{Pr}(k,\widehat{\theta })`$ and calculate $`D^{}`$ and $`D^{}`$, respectively, for each replicate. From $`L`$ bootstrap-replicates we obtain the Null distribution of $`D`$ and $`D^{}`$ under the estimated degree distribution. For sufficiently large $`L`$ we can thus determine approximate $`p`$-values which allow us to test if the empirical degree distribution is commensurate with the estimated degree distribution.
## 3 Statistical analysis of biological networks
Here we apply the analysis of the preceeding sections to the study of PINs and metabolic networks. It is easy to find a straight-line fit to some degree interval for all of the datasets considered here. For a powerlaw to be meaningful it has to extend over at least two or three decades and with a maximum degree of $`k_{\mathrm{max}}300`$ this will be unachievable for the present data sets. We therefore use all the data and fit the model which yields the best overall a description of the degree distribution.
### 3.1 Analysis of PIN data
In table 2 we show the maximum composite likelihoods for the degree distributions calculated from PIN data collected in five model organisms xenarios\_NAR2000 (the protein interaction data was taken from the DIP data-base; http://dip.doe-mbi.ucla.edu). We find that the standard scale-free model (or its finite size versions) never provides the best fit to the data; in three networks (C.elegans, S.cerevisiae and E.coli) the lognormal distribution (M5) explains the data best. In the remaining two organisms the stretched exponential model provides the best fit to the data. The bold likelihoods correspond to the highest Akaike weights. Apart from the case of H.pylori (where $`\mathrm{max}(w_j)=w_50.95`$ for M5 and $`w_60.05`$) the value of the maximum Akaike weight is always $`>0.9999`$. For C elegans, however, the scale-free model and its finite size versions are better than the lognormal model, M5.
For the yeast PIN the best fit curves (obtained from the MLEs of the parameters of models M1-M6) are shown in figure 1, together with the real data. Visually, log-normal (green) and stretched exponential (blue) appear to describe the date almost equally well. Closer inspection, guided by the Akaike weights, however, shows that the fit of the lognormal to the data is in fact markedly better than the fit of the stretched exponential. But the failure of quickly decaying distributions such as the Poisson distribution, characteristic for classical random graphs bollobas\_book1 to capture the behaviour of the PIN degree distribution is obvious.
Interestingly, common heuristic finite size corrections to the standard scale-free model improve the fit to the data (measured by the AIC). But compared to the lognormal and stretched exponential models they still fall short in describing the PIN data in the five organisms.
Figure 2 shows only the three curves with the highest values of $`\omega _j`$, which apart from E.coli are the log-normal, stretched exponential and power-law distributions; for E.coli, however, the Gamma distribution replaces the power-law distribution. These figures show that, apart from C.elegans the shape of the whole degree distribution is not power-law like, or scale-free like, in a strict sense. Again we find that log-normal and stretched exponential distributions are hard to distinguish based on visual assessment alone. Figures 1 and 2, together with the results of table 2, reinforce the well known point that it is hard to choose the best fitting function based on visual inspection. It is perhaps worth noting, that the PIN data is more complete for S.cerevisiae and D.melanogaster than for the other organisms.
The standard scale-free model is superior to the log-normal only for C.elegans. The order of models (measured by decreasing Akaike weights) is M6, M5, M4, M2, M3, M1 for D.melanogaster, M6, M4, M5, M2, M3, M1 for C.elegans, M5, M6, M4, M2, M3, M1 for S.cerevisiae and H.pylori, and M5, M3, M6, M4, M2, M1 for E.coli. Thus in the light of present data the PIN degree distribution of E.coli lends more support to a Gamma distribution than to a scale-free (or even stretched scale-free) model. There is of course, no mechanistic reason why the gamma distribution should be biologically plausible but this point demonstrates that present PIN data is more complicated than predicted by simple models. Therefore statistical model selection is needed to determine the extent to which simple models really provide insights into the intricate architecture of PINs. For completeness we note that model selection based on BIC results in the same ordering of models as the AIC shown here.
In table 3 we give the values of $`D`$ and $`D^{}`$ for the empirical degree distributions. The estimated cumulative distribution $`P(k;\widehat{\theta })`$ was obtained from the maximum likelihood fits of the respective models. The results in table 3 show that the maximum likelihood framework (or the respective models) sometimes cannot adequately describe the tails of the distribution —at low and high values of the connectivity— in some cases, where $`D^{}>>D`$. The order of the different models suggested by $`D`$ for the three models generally agrees with the ordering obtained from the AIC.
### 3.2 Analysis of MN data
Metabolic networks aim to describe the biochemical machinery underlying cellular processes. Here we have used data from the KEGG database (www.genome.jp/kegg) with additional information from the BRENDA database (www.brenda.uni-koeln.de). The nodes are the enzymes and an edge is added between two enzymes if one uses the products of the other as educts.
From figure LABEL:metabfig and table 4 it is apparent that the maximum likelihood scale-free model obtained from the whole network data does not provide an adequate description of the MN data. This should, however, not be too surprising as the network is only relatively small with maximum degree $`k_{\mathrm{max}}=83`$. The degree distribution appears to decay in an essentially exponential fashion but the stretched exponential has the required extra flexibility to describe the whole of the data better than the other models. Measured by goodness of fit statistics $`D`$ and $`D^{}`$, however, the log-normal model ($`D=0.02`$ and $`D^{}=0.06`$) performs rather better than the stretched exponential ($`D=0.07`$ and $`D=0.22`$); the scale-free model again performs very badly ($`D=1.0`$ and $`D^{}=34.2`$) for the metabolic network data.
## 4 Conclusions
We have shown that it is possible to use standard statistical methods in order to determine which probability model describes the degree distribution best. We find that the common practice of fitting a pure power-law to such experimental network datador\_book ; albert\_RMP2002 may obscure information contained in the degree distribution of biological networks. This is often done by identifying a range of connectivities from the log-log plots of the degree distribution which can then be fitted by a straight line. Not only is this wasteful in the sense that not all of the data is used but it may obfuscate real, especially finite-size, trends. The same will very likely hold true for other biological networks, too may\_PRE2001 . The approach used here, on the other hand, (i) uses all the data, and (ii) can be extended to assessing levels of confidence through combining a bootstrap procedure with the Akaike weights. What we have shown here, in summary, is that statistical methods can be usefully applied to protein interaction and metabolic network data.
Even though the degree distribution does not capture all (or even most) of the characteristics of real biological networks there is reason to reevaluate previous studies. We find that formally real biological networks provide very little support for the common notion that biological networks are scale-free. Other fat-tailed probability distributions provide qualitatively and quantitatively better descriptions of the degree distributions of biological networks. For protein interaction networks we found that the log-normal and the stretched exponential offer superior descriptions of the degree distribution than the powerlaw or its finite size versions. For metabolic versions our results confirms this. Even the exponential model outperformed the scale-free model in describing the empirical degree distribution (we note that randomly growing networks are characterized by an exponentially decreasing degree distribution). The best models are all fat-tailed —like the scale-free models— but are not formally scale-free. Unfortunately, there is as yet no known physical model for network growth processes that would give rise to log-normal or stretched exponential degree distributions.
There is thus a need to develop and study theoretical models of network growth that are better able to describe the structure of existing networks. This probably needs to be done in light of at least three constraints: (i) real networks are finite sized and thus, in the terms of statistical physics, mesoscopic systems; (ii) present network data are really only samples from much larger networks as not all proteins are included in present experimental setup (in our case S.cerevisiae has the highest fraction, 4773 out of approximately 5500-6000 proteins); the sampling properties have recently been studied and it was found that generally the degree distribution of a subnet will differ from that of the whole network. This is particularly true for scale-free networks. (iii) biological networks are under a number of functional and evolutionary constraints and proteins are more likely to interact with proteins in the same cellular compartment or those involved in the same biological process. This modularity —and the information already availabe e.g. in gene ontologies— needs to be considered. Finally there is an additional caveat: biological networks are not static but likely to change during development. More dynamic structures may be required to deal with this type of problem.
Quite generally we believe that we are now at a stage where simple models do not necessarily describe the data collected from complex processes to the extent that we would like them to. But as Burda, Diaz-Correia and Krzywicki point out burda\_PRE2001 , even if a mechanistic model is not correct in detail, a corresponding statistical ensemble may nevertheless offer important insights. We believe that the statistical models employed here will also be useful in helping to identify more realistic ensembles.
The maximum likelihood, goodness of fit and other tools and methods for the analysis of network data are implemented in the NetZ R-package which is available from the corresponding author on request.
Acknowledgements: We thank the Wellcome Trust for a research fellowship (MPHS) and a research studentship (PJI). CW is supported by the Danish Cancer Society. Financial support from the Royal Society and the Carlsberg Foundation (to MPHS and CW) is also gratefully acknowledged. We have furthermore benefitted from discussions with Eric de Silva, Bob May and Mike Sternberg.
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# Anharmonic oscillator and double-well potential: approximating eigenfunctions
México ICN-UNAM 05-03
June 1, 2005
Abstract
> A simple uniform approximation of the logarithmic derivative of the ground state eigenfunction for both the quantum-mechanical anharmonic oscillator and the double-well potential given by $`V=m^2x^2+gx^4`$ at arbitrary $`g0`$ for $`m^2>0`$ and $`m^2<0`$, respectively, is presented. It is shown that if this approximation is taken as unperturbed problem it leads to an extremely fast convergent perturbation theory.
Invited contribution to Letters in Mathematical Physics
to a Special Issue in memory of Professor Felix A. Berezin
For the last fifty years the quantum one-dimensional anharmonic oscillator $`(m^20)`$ as well as the double-well potential problem $`(m^2<0)`$ described by the Hamiltonian
$$=\frac{d^2}{dx^2}+m^2x^2+gx^4,$$
(1)
permanently attracted a lot of attention. The interest to these problems ranges from various branches of physics to chemistry and biology. It can not be an exaggeration to say that after the seminal papers by C. Bender-T.T. Wu at 1969-1973 BW in near thousand of physics articles the problem (1) was touched in one way or another. This seemingly simple problem revealed extremely rich internal structure which looks intrinsic for any non-trivial problem of quantum mechanics and even for quantum field theory. In particular, the obvious failure of the perturbation theory in powers of $`g`$ due to its asymptotic (divergent) nature in (1) pushed the development of non-perturbative methods. In practice, the anharmonic oscillator (1) served always as a test-ground for non-perturbative methods. Another important feature of the anharmonic oscillator (1) is related to the fact that it can be interpreted as one-dimensional quantum field theory $`g\varphi ^4`$ with zero spacial dimension, $`(1,0)`$. Hence, a study of (1) can be insightful to a realistic four-dimensional quantum field theory $`g\varphi ^4`$. The present article is devoted to the construction of a uniform approximation of the ground state eigenfunction of (1) in $`x`$-space which would continue to be valid for different values of the coupling constant $`g`$ and the parameter $`m^2`$. If such an approximation is constructed the energy of the ground state as well as different average values can be calculated with a guaranteed accuracy. We are not aware of previous attempts to construct a uniform approximation of the eigenfunctions.
Take the Schroedinger equation for (1)
$$\frac{d^2\mathrm{\Psi }}{dx^2}+m^2x^2\mathrm{\Psi }+gx^4\mathrm{\Psi }=E(m^2,g)\mathrm{\Psi },_{\mathrm{}}^+\mathrm{}\left|\mathrm{\Psi }\right|^2𝑑x<\mathrm{}.$$
(2)
For $`m^20`$ the potential in (2) is a single-well potential, which describes the celebrated quartic anharmonic oscillator. For $`m^2<0`$ the potential in (2) is a two-well potential. It describes another celebrated potential - the so-called double-well potential (also known as the Higgs potential, Lifschitz potential). It is easy to check that for the energy in (2) the Symanzik scaling relation holds
$$E(m^2,g)=g^{1/3}E(\frac{m^2}{g^{2/3}},1),\mathrm{\Psi }(x;m^2,g)=\mathrm{\Psi }(xg^{1/6};\frac{m^2}{g^{2/3}},1).$$
This manifests that the original problem (2) is essentially a single-parametric problem.
Eigenfunctions of (2) are sharply changing functions in $`x`$ being characterized by a power-like behavior at $`|x|0`$ and an exponentially-decaying one at $`|x|\mathrm{}`$. In order to make them smooth we introduce the exponential representation for eigenfunctions
$$\mathrm{\Psi }\left(x\right)=e^{\phi \left(x\right)}.$$
(3)
Following the oscillation (Sturm) theorem the function $`\phi `$ should have logarithmic singularities at real $`x`$ which correspond to nodes of the wavefunction. Its regular part which appears after subtraction of those singularities should be slow-changing function of the power-like behavior at both small and large distances. After substitution of (3) into (2) we get a Riccati equation
$$y^{}y^2=Em^2x^2gx^4,y=\phi ^{}=\left(\mathrm{log}\mathrm{\Psi }\left(x\right)\right)^{},y\left(0\right)=0.$$
(4)
This famous equation serves as a basis for developing the WKB approximation scheme. Due to symmetric nature of the r.h.s. of (4) $`xx`$ the function $`y`$ is odd. Its structure for the $`n`$th excited state should be
$$y=\underset{i=1}{\overset{n}{}}\frac{1}{xx_i}+y_{reg}\left(x\right),$$
(5)
where $`x_i`$ is the position of $`i`$th node. The regular part $`y_{reg}(x)`$ should be a non-singular at real $`x`$ odd function which vanish at $`x=x_i,i=1,2,\mathrm{}n`$. It is evident that all non-zero $`x_i`$ come in pairs: once we have a node at $`x=x_i`$ always there exists a node at $`x=x_i`$.
Now we consider the case of the ground state. According to above, in the phase (5) $`n=0`$ and singular part of $`y`$ is absent, $`y=y_{reg}`$. Hence, the function $`y`$ has no singularities at real $`x`$. It is easy to find its asymptotic behavior
$$y=g^{1/2}x\left|x\right|+\frac{m^2}{2g^{1/2}}\frac{\left|x\right|}{x}+\frac{1}{x}\frac{4gE+m^4}{8g^{3/2}}\frac{1}{x\left|x\right|}\frac{m^2}{2g}\frac{1}{x^3}+\mathrm{}\text{at}\left|x\right|\mathrm{},$$
(6)
while
$$y=Ex+\frac{E^2m^2}{3}x^3+\frac{2E\left(E^2m^2\right)3g}{15}x^5+\mathrm{}\text{at}\left|x\right|0.$$
(7)
One can demonstrate that for any excited state the regular part $`y_{reg}`$ has similar formal expansions at $`|x|0,\mathrm{}`$.
Let us develop a certain iterative procedure (perturbation theory) for solving the Riccati equation (4). From the point of finding the wave function it will be a multiplicative perturbation theory unlike a standard additive perturbation theory. Such a perturbation theory was developed for the first time by Price Price and then it was rediscovered many times. Eventually, it was called the ‘Logarithmic Perturbation Theory’ (for the history remarks and discussion see Turbiner:1984 and references therein).
As a first step we choose some square-integrable function $`\mathrm{\Psi }_0`$ and calculate its logarithmic derivative
$$y_0=\left(\mathrm{log}\mathrm{\Psi }_0\right)^{}=\frac{\mathrm{\Psi }_0^{}}{\mathrm{\Psi }_0}.$$
(8)
It is clear that $`\mathrm{\Psi }_0`$ is the exact eigenfunction of the Schroedinger operator with a potential
$$V_0=\frac{\mathrm{\Psi }_0^{\prime \prime }}{\mathrm{\Psi }_0}=y_0^2y_0^{},$$
(9)
where without a loss of generality we put their eigenvalue equals to zero, $`E_0=0`$. It is nothing but a choice of the reference point for eigenvalues. Now we can construct a perturbation theory for Riccati equation taking $`\mathrm{\Psi }_0`$ and $`y_0,V_0`$ as zero approximation, which characterizes the unperturbed problem. One can write the original potential $`V=m^2x^2+gx^4`$ as a sum,
$$V=V_0+\left(VV_0\right)V_0+V_1,$$
(10)
thus, taking a deviation of the original potential from the potential of the zero approximation as a perturbation. We always can insert a formal parameter $`\lambda `$ in front of $`V_1`$ and develop a perturbation theory in powers of $`\lambda `$,
$$E=\underset{k=0}{\overset{\mathrm{}}{}}\lambda ^kE_k,y=\underset{k=0}{\overset{\mathrm{}}{}}\lambda ^ky_k,$$
(11)
putting $`\lambda =1`$ afterwards. Perhaps, it is worth emphasizing that in spite of the fact that we study iteratively the equation (4), in general, this perturbation series has nothing to do with a standard WKB expansion. By substituting (11) into (4) we arrive at the equations which defines iteratively the corrections
$$y_k^{}2y_0y_k=E_kQ_k,$$
(12)
where
$`Q_1`$ $`=`$ $`V_1,`$
$`Q_k`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{k1}{}}}y_iy_{ki},k=2,3,\mathrm{}`$
It is interesting that the operator in the l.h.s. of (12) does not depend on $`k`$, while $`Q_k`$ in the r.h.s. can be interpreted as a perturbation on the level $`k`$. The solution of (12) can be found explicitly and is given by
$`E_k`$ $`=`$ $`{\displaystyle \frac{_{\mathrm{}}^{\mathrm{}}Q_k\mathrm{\Psi }_0^2𝑑x}{_{\mathrm{}}^{\mathrm{}}\mathrm{\Psi }_0^2𝑑x}},`$ (13)
$`y_k`$ $`=`$ $`\mathrm{\Psi }_0^2{\displaystyle _{\mathrm{}}^x}\left(E_kQ_k\right)\mathrm{\Psi }_0^2𝑑x^{}.`$ (14)
It is easy to demonstrate that
$$\left|y_1\right|\text{Const},$$
(15)
provides a sufficient condition for this perturbation theory (11) to be convergent Turbiner:1984 . Note that this condition is very rough and can be strengthened.
The first two terms in the expansion of energy (11) in the above-described perturbation theory admit an interpretation in the framework of the variational calculus. Let us assume that our variational trial function $`\mathrm{\Psi }_0(x)`$ is normalized to 1. We can calculate the potential $`V_0`$ where $`\mathrm{\Psi }_0(x)`$ is the ground state eigenfunction and even put $`E_0=0`$ (see a discussion above). Pure formally, we construct the Hamiltonian $`H_0=p^2+V_0`$ for which $`H_0\mathrm{\Psi }_0(x)=0`$. The variational energy is equal to
$`E_{var}`$ $`=`$ $`{\displaystyle \psi _0H\psi _0}=\underset{=E_0}{\underset{}{{\displaystyle \psi _0H_0\psi _0}}}+\underset{=E_1}{\underset{}{{\displaystyle \psi _0\underset{VV_0}{\underset{}{\left(HH_0\right)}}\psi _0}}}`$ (16)
$`=`$ $`E_0+E_1\left(V_1=VV_0\right)E_{exact}.`$
Of course, $`\mathrm{\Psi }_0(x)`$ could depend on free parameters. In this case both $`V_0`$ and $`V_1`$ depend on parameters as well. Minimization of $`E_{var}`$ with respect to the parameters can be performed and the variational principle guarantees that $`E_{var}`$ gives upper bound to the ground state energy. This simple interpretation (16) reveals a fundamental difference between perturbation theory and variational calculus. Variational estimates can be obtained independently on the fact that the perturbation theory associated with $`\mathrm{\Psi }_0(x)`$ is convergent or divergent. However, it seems natural to remove this difference by requiring a convergence of the perturbation series. In this case by calculating the next terms $`E_2,E_3,\mathrm{}`$ in (11) one can estimate the accuracy of variational calculation from one side and improve it iteratively from another side. An immediate criteria how to choose $`\mathrm{\Psi }_0(x)`$ in order to get a convergent perturbation theory is to have the perturbation potential $`V_1`$ to be subordinate with respect to the non-vanishing potential of zero approximation $`V_0`$,
$$\left|\frac{V_1}{V_0}\right|<1,\text{for}\left|x\right|>R.$$
(17)
This requirement has extremely non-trivial physical implication: in order to guarantee a convergence of perturbation theory a domain where the wavefunction is exponentially small should be reproduced as precise as possible. The same time a description of a domain where the wavefunction is of the order 1 is not important. It contradicts to a straightforward physics intuition and underlying idea of variational calculus which, in particular, requires a precise description of the domain where the wavefunction is of the order 1. Needless to say that namely the latter domain gives a dominant contribution to the integrals which define the energy in the variational calculations.
In our approach the main object for study is the logarithmic derivative of the wavefunction $`y`$ (see (4)). Since $`y`$ is antisymmetric, we will construct different interpolations of $`y`$ between $`x=0`$ and $`x=\mathrm{}`$ (see (6),(7)). In order to fulfil the requirement of the convergence (17) it is enough to take into account in these interpolations the leading term of the asymptotics at $`|x|\mathrm{}`$. Then the interpolation we have built is taken as zero approximation in our perturbation theory (11).
1. The simplest interpolation can be written as follows
$$y_0=ax+b\sqrt{g}x\left|x\right|,$$
(18)
where $`a,b`$ are parameters. These parameters can be fixed either by taking them as variational and making a minimization of (16), $`a_{min},b_{min}`$, or following the idea to reproduce exactly the leading asymptotic behavior of $`y`$ at $`x=\mathrm{}`$ (see (6)), which requires to put $`b=1`$.
The ground state eigenfunction which corresponds to (18)
$$\psi _0^{\left(1\right)}=\mathrm{exp}\left\{\frac{ax^2}{2}b\frac{\sqrt{g}}{3}\left|x\right|^3\right\},$$
(19)
is the exact one for the potential
$$V_0=a^2x^2+b^2gx^42b\sqrt{g}\left|x\right|\left(1ax^2\right),E_0=a.$$
The perturbation potential $`V_1=VV_0`$ is of the form
$$V_1=\left(m^2a^2\right)x^2+\left(1b^2\right)gx^4+2b\sqrt{g}\left|x\right|\left(1ax^2\right),$$
where the first two terms would disappear if we place $`a=\pm m`$, and $`b=1`$. In this case the first two terms in the perturbation theory are
$$E^{\left(1\right)}E_0+E_1=m+2\sqrt{g}\frac{_0^{\mathrm{}}x\left(1mx^2\right)e^{mx^2\frac{2\sqrt{g}}{3}x^3}𝑑x}{_0^{\mathrm{}}e^{mx^2\frac{2\sqrt{g}}{3}x^3}𝑑x}.$$
(20)
This expression leads to sufficiently high relative accuracy $`10^2`$ (comparing to the accurate numerical results) for any $`g0`$ and value of $`m^2`$. It can be confirmed by calculation of the second correction $`E_2`$. However, a slight modification of of the interpolation (18) by including a term $`1/x`$, which is presented in (6), in a form
$$\stackrel{~}{\psi }_0^{\left(1\right)}=\frac{1}{\sqrt{1+cx^2}}\mathrm{exp}\left\{\frac{ax^2}{2}b\frac{\sqrt{g}}{3}\left|x\right|^3\right\},$$
(21)
where $`c`$ is a variational parameter, immediately leads to a drastic increase in accuracy (see Table 1). According to a quite simple, straightforward analysis of the second correction $`E_2`$ the variational energy deviates from exact one in $`10^3`$ in relative units for both anharmonic oscillator and double-well potential for studied values of the parameters in (1) (see Table 1). When the second correction is taken into account the relative deviation reduces in two orders of magnitude becoming $`10^5`$. It is worth mentioning that the quality of the approximation is reflected in the fact that $`b_{min}`$ deviates from the exact value $`b=1`$ in several percent.
The ground state function is symmetric w.r.t. $`xx`$. Hence, it has an extremum at $`x=0`$. For any fixed $`m^2<0`$ there exists a value $`g_{crit}`$ such that for $`g>g_{crit}`$ this extremum is a maximum, otherwise a minimum. It is easy to find out that the critical point $`g=g_{crit}`$ corresponds to the vanishing ground state energy, $`E=0`$. Using the function (21) it was calculated the critical value $`g_{crit}=0.302405`$ for $`m^2=1`$. It is quite interesting from physical point of view that for a family of double-well potentials with fixed $`g`$ there exists a domain $`0>m^2>(m^2)_{crit}`$ where the ground-state eigenfunction has the maximum at the origin, which corresponds to the position of the unstable equilibrium similar to what takes place for the single-well case. For example, if $`g=1`$, the value of $`(m^2)_{crit}=2.219597`$. It implies that the particle in such a potential with the ground state energy above the barrier, $`E>0`$, somehow does not feel the existence of two minima.
2. The expression $`E^{(1)}`$ (20) reproduces the harmonic oscillator energy in the limit $`g0`$. At $`g=0`$ the energy $`E^{(1)}`$ has a singularity as it must be. However, the expansion of (20) in powers of $`g`$ contains besides the integer powers in $`g`$ also half-integer powers which must not be present in the formal expansion of the energy in $`g`$ (see (11). Also the interpolation (18) reproduces the leading terms only in the expansions at $`|x|0,\mathrm{}`$. It is a definite drawback of the interpolation (18) as well as (21) and it should be fixed. One of the simplest ways to fix it is to take the function
$$\psi _0^{(0,0)}=\frac{1}{\sqrt{1+c^2x^2}}\mathrm{exp}\left\{\frac{A+ax^2/2+bgx^4/4}{\left(d^2+gx^2\right)^{1/2}}\right\},$$
(22)
where $`A,a,b,c,d`$ are variational parameters. If the parameters are chosen to be
$$b=\frac{4}{3},a=\frac{d^2}{3}+m^2,$$
(23)
the dominant and the first two subdominant terms in the asymptotic expansion (6) are reproduced by the function (22) exactly. In this case the convergency condition (15) is satisfied, the data for $`|y_1|_{max}`$ are in Table II. These two cases we call the Case 1 (two parameters are fixed (see (23)) to reproduce the growing terms in asymptotic expansion of $`y`$) and Case 2 (all five parameters in (22) are variational), respectively. Table II demonstrates that even the variational energy $`E^{(1)}`$ already provides extremely high accuracy: $`10^810^9`$ (Case 1) and $`10^{10}10^{11}`$ (Case 2) for studied values of $`m^2,g`$. These results reproduce (or exceed) the best known numerical results in literature. The order of the third energy correction $`E_3`$ is already in the three-four orders of magnitude less than the second correction $`E_2`$. It indicates to the extremely fast convergence of the series in the parameter $`\lambda `$ (see (11)). The energies $`E^{(3)}(=E^{(1)}+E_2+E_3)`$ for given values of $`m^2,g`$ are the most accurate among known in literature for the moment providing 16-17 significant digits (not all these digits are shown in Table II).
In the Case 1 we are focused to reproduce the behavior of the wavefunction at large distances while the behavior at small distances is defined as a result of the variational procedure. We can calculate the coefficient in front of $`x`$ in the logarithmic derivative of (22) which we denote $`E_{exp}`$ (see Table II). It should be the exact energy of the ground state in the case of the exact $`y`$ (see (7)). It is worth emphasizing that the deviation of $`E_{exp}`$ from $`E^{(1)}`$ (or from the exact energy $`E`$) appears in the fifth significant digit (!).
The important question concerns to the deviation of the function (22) from the exact function. We consider the Case 1 when the growing asymptotics of $`y`$ at large distances is reproduced exactly. Due to this fact the perturbation theory for any studied $`m^2,g`$ is fast convergent and $`y_1`$ is bounded for any real $`x`$,
$$\left|y_1\right|_{max}0.01.$$
Hence, the function (22) provides the uniform approximation of the exact $`y`$. In Fig.1 for $`m^2=1,g=2`$ the behavior of $`y_0`$ it is shown while Fig.2 contains the behavior of the first correction $`y_1`$. It is worth mentioning that the maximal value of $`y_1`$: $`|y_1|_{max}0.0029`$ appears at $`x3.9`$ where the value of $`\psi _0^{(0,0)}`$ (22) is already extremely small, $`10^9`$ and $`y_015.2`$. Then $`y_11/x^2`$ at $`x1`$. The ratio $`y_1/y_0`$ (see Fig.3) is the extremely small function which tends to zero at $`|x|\mathrm{}`$. A similar situation appears for other values of $`m^2,g`$. For the double-well potential while approaching the semiclassical limit: $`m^2\mathrm{}`$ and $`g`$ is kept fixed, e.g. $`g=1`$ \- the accuracy provided by (22) remains unchanged: $`E_21.\times 10^6`$ for $`m^2`$ ranging from -1 to -30. In the same range the parameters of (22) are varied very little, being always of the same order of magnitude, although the value of the energy changes from $`1`$ to $`105`$. It was also calculated the parameters of the potential (1) which the ground state energy vanishes,
$$E\left(m^2=2.2195970861,g=1\right)10^{13}.$$
It is worth mentioning that uniform approximation for eigenfunction of $`(2k+p)`$th excited state, $`k=0,1,2,\mathrm{},p=0,1`$ can be easily constructed as well. It has a form
$$\psi _0^{(k,p)}=\frac{x^pP_k\left(x^2\right)}{\sqrt{1+c^2x^2}}\mathrm{exp}\left\{\frac{A+ax^2/2+bgx^4/4}{\left(d^2+gx^2\right)^{1/2}}\right\},$$
(24)
where $`P_k`$ is a polynomial of $`k`$th degree with positive roots. Parameters $`a,b`$ are always chosen following (23) in order to reproduce the growing terms in the asymptotic expansion of $`y`$ (Case 1 for the ground state)). In order to find the other parameters the orthogonality constraint should be imposed: for fixed $`(k,p)`$, (24) should be orthogonal to previously constructed functions $`(m,p)`$, at $`m=0,1,2,\mathrm{},k1`$. It fixes some parameters in (24) while the remaining parameters are used as variational. It is worth mentioning that very important physical characteristic of (1) at $`m^2<0`$ is the energy gap between the energies of the ground state $`(0,0`$ and the first excited state $`(0,1`$ in the semiclassical limit: $`m^2`$ is kept fixed and $`g\mathrm{}`$. Although we do not provide concrete numerical results in this paper, this limit is reproduced with very high accuracy.
As a conclusion I would like to recollect my first and, in fact, the only scientific encounter with F.A. Berezin. It was in mid-1970s just before his breakthrough results in supermathematics. I was a recent graduate of Moscow Institute for Physics and Technology just hired by Institute for Theoretical and Experimental Physics. Among theoretical physicists F.A. Berezin had a reputation of the unique mathematician who was able to understand what physicists are talking about and with whom one could discuss your mathematical difficulties. He had run a weekly seminar at Mathematics Department of Moscow State University. By accident I went to one of his seminars but I could not say that I understood much. During the seminar break in the moment when I was ready to leave he approached to me, introduced himself and asked what I am working on. I answered that for already several years I had been trying desperately to find an analytic solution of the anharmonic oscillator $`x^4`$ (see (1)). Or, at least, to understand why all my attempts had failed. He immediately said: ”I know an anharmonic oscillator, which has an analytic solution! Take a function $`e^{x^4}`$, differentiate it twice and divide the result by $`e^{x^4}`$. This is the anharmonic oscillator potential where you know an exact eigenfunction.” Then we talked for several minutes. In the end, he said that together with M.A. Shubin (now a Distinguished Professor at Northeastern University in Boston) they would write a book on the Schroedinger equation and one day if I did not mind they could ask my opinion or advice. He introduced M.A. Shubin to me with whom I become a lifelong friend. After the seminar F.A. Berezin, M.A. Shubin and myself took a walk to the nearest metro station. I remember it was a very interesting conversation about mathematics but I do not remember what it was about - I was engrossed with thinking about the remarks of F.A. Berezin.
This short meeting had struck me and influenced strongly my scientific life. From human point of view I was deeply impressed how serious and friendly he was towards me, almost still a student. After some time I understood that instead of solving the Schroedinger equation with an already given potential one can generate a zillion potentials for which a single eigenstate is known exactly. It led me to an idea of choosing one of such potentials as a zero approximation to construct a convergent iterative procedure for a given eigenstate, not for the whole spectra Turbiner:1979 . Another idea was related to a natural question of how to construct a potential for which not one but two or more eigenstates can be found explicitly. It led to a discovery of quasi-exactly-solvable quantal problems Turbiner:1988 . A few years after our meeting F.A. Berezin tragically died and the chance to discuss with him disappeared, but all my life I keep a memory of him. Recently, I told this story to one of the greatest minds in today’s mathematics who was among very close colleagues and friends of F.A. Berezin. He responded with sorrow that ”… only after his death we realized how strong was his influence on all of us and, eventually, how great he was!”
Acknowledgments
Author thanks J.C. Lopez Vieyra the interest to the work and a help with computer calculations. The work is supported in part by DGAPA grant No.IN124202 (Mexico).
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# The vacuum state functional of interacting string field theory
## 1 Introduction and review
The Euclidean vacuum functional of quantum field theory and first quantised string theory can be constructed via a large time functional integral. The Lorentzian continuation in field theory is
$$\mathrm{\Psi }_0[\varphi ]=𝒟\phi \mathrm{exp}\left(\frac{i}{\mathrm{}}\underset{\mathrm{}}{\overset{0}{}}dx^0L(\phi ,\dot{\phi })\right)|_.^{\phi =\varphi \text{ at time }0}$$
(1.1)
The boundary condition in the infinite past is that the field must be regular. By shifting the integration variable by a piece proportional to the Heaviside function the field dependence is moved out of the boundary conditions and into the action, giving
$$\mathrm{\Psi }_0[\varphi ]=𝒟\phi \mathrm{exp}(\frac{i}{\mathrm{}}\underset{\mathrm{}}{\overset{0}{}}\mathrm{d}x^0L(\phi ,\dot{\phi })+\frac{i}{\mathrm{}}\mathrm{d}^D𝐱\dot{\phi }(𝐱,0)\varphi (𝐱))|^{\phi =0\text{ at time }0}.$$
(1.2)
Standard field theory results then imply that the logarithm of the vacuum functional is the sum of connected Feynman diagrams constructed from a propagator $`G_d`$ which obeys Dirichlet boundary conditions on the hypersurface $`x^0=0`$ with vertices integrated over the interval $`x^0<0`$. All external legs (those attached to $`\varphi `$) end on the boundary with a time derivative which results from $`\varphi `$ being coupled to $`\dot{\phi }`$.
Our aim is to construct the vacuum state functional for string field theory, following the interest generated by the Sen conjectures . The string field interaction will be non-local in our chosen time co-ordinate, but the above description of the vacuum fails if the interaction becomes non-local. This is because we would attempt to describe the functional by a sum over field histories on the half space $`x^0<0`$, but through a non-local interaction the field can couple to itself at arbitrary times.
Following this introduction we will describe an alternative method of constructing the vacuum functional which applies to both local and non-locally interacting field theories, and then generalise to string field theory. The essential ingredient is the set of gluing properties. For the scalar field propagator these are
$$\mathrm{d}^D𝐲G_0(𝐱_2,t_2;𝐲,t)\left(2\frac{}{t}\right)G_0(𝐲,t;𝐱_1,t_1)=\{\begin{array}{cc}iG_0(𝐱_2,t_2;𝐱_1,t_1)& \hfill t_2>t>t_1\\ iG_0(𝐱_2,t_2;𝐱_1,t_1)& \hfill t_1>t>t_2\\ iG_I(𝐱_2,t_2;𝐱_1,t_1)& \hfill t>t_1,t_2\\ iG_I(𝐱_2,t_2;𝐱_1,t_1)& \hfill t<t_1,t_2\end{array}$$
(1.3)
These rules apply to both the open and closed string field propagator when time is associated with the hypersurfaces $`X^0(\sigma )=\text{ constant}`$. The point $`𝐱^i`$ is replaced by the 25 dimensional spacetime curve $`𝐗^i(\sigma )`$ along with the non-vanishing ghost components on the Dirichlet sections of the propagator. In this paper $`𝐗`$ will be shorthand for this whole set to simplify this notation. For a proof of the gluing rules and a discussion of the ghost sector see .
We will construct the vacuum wave functional using the vacuum expectation values it must generate. This will apply to both open and closed string field theories with a cubic interaction, though the explicit form of this interaction is arbitrary. Our construction is perturbative, we give the first few terms of the functional explicitly, but the procedure used generalises to all orders in perturbation theory. In the final section of this paper we will return to local field theory and use our gluing rules to give a diagrammatic demonstration of the equivalence of Schrödinger and covariant field theory pictures in field theory.
## 2 The string field vacuum state
The vacuum state functional is the generator of vacuum expectation values. For example, the open string three field expectation value would be represented by, to leading order in the coupling, the diagram
(2.1)
Such expectation values can be used to reconstruct the vacuum state functional. We will demonstrate this for quantum field theory and then generalise to strings. Since the string field interaction is in various guises cubic we will consider a non-local scalar $`\varphi ^3`$ theory with interaction Hamiltonian density
$$\frac{\lambda }{3!}\underset{j=1}{\overset{3}{}}\mathrm{d}^D𝐱_j\mathrm{d}t_jW(t,𝐱;𝐱_1,t_1,𝐱_2,t_2,𝐱_3,t_3)\varphi (𝐱_1,t_1)\varphi (𝐱_2,t_2)\varphi (𝐱_3,t_3)$$
(2.2)
written in terms of some kernel $`W`$. Replacing this kernel with a product of delta functions $`\delta (tt_i)`$ recovers the local $`\varphi ^3`$ Hamiltonian density. The technique we are about to describe applies to both local and non-local interactions, as demonstrated in for local $`\varphi ^4`$ theory.
To begin, consider the free field vacuum functional, which must yield the equal time propagator as a vacuum expectation value,
$$\mathrm{}G(𝐱,0;𝐲,0)=\mathrm{\hspace{0.17em}0}|\varphi (𝐱,0)\varphi (𝐲,0)|\mathrm{\hspace{0.17em}0}=𝒟\varphi \varphi (𝐱)\varphi (𝐲)\mathrm{\Psi }_0^2[\varphi ]$$
which implies that the free field vacuum is the exponent of the inverse of the equal time propagator. This inverse is
$$G_0(𝐱,0;𝐲_0)^1=4\frac{^2}{tt^{}}G_0(𝐱,t;𝐲_t^{})|_{t=t^{}=0}$$
(2.3)
as we proved in , so the free field vacuum is
$$\mathrm{\Psi }_0^{\text{free}}[\varphi ]=\mathrm{exp}\left(\frac{1}{4\mathrm{}}\mathrm{d}^D(𝐱,𝐲)\varphi (𝐱)G_0(𝐱,\stackrel{}{0};𝐲,\stackrel{}{0})\varphi (𝐲)\right),$$
(2.4)
as is easily checked. The bullet is a time derivative with factor $`2`$, as in (1.3). Now introduce the cubic interaction. Our arguments are based in the Schrödinger representation, and treats all spatial arguments in the same way as covariant methods. Therefore we may choose the interaction to be local or non-local in space without affecting the results, as such we will write the Hamiltonian as
$$\frac{\lambda }{3!}\underset{j=1}{\overset{3}{}}\mathrm{d}t_jW(t;t_1,t_2,t_3)\varphi (t_1)\varphi (t_2)\varphi (t_3)$$
(2.5)
to keep notation compact, suppressing the spatial dependencies whenever possible. Accordingly we will abbreviate our representation of the free propagator to
$$\begin{array}{cc}\hfill G_𝐱(0;t_j)& :=G_0(𝐱,0;𝐱_j,t_j),\hfill \\ \hfill \varphi _𝐱G_𝐱(0;t_j)& :=\mathrm{d}^D𝐱\varphi (𝐱)G_0(𝐱,0;𝐱_j,t_j).\hfill \end{array}$$
(2.6)
To clarify our arguments we also assume that the interaction kernel is symmetric in its indices. The following applies when this is not the case, the only difference being the appearance of sums of terms corresponding to inequivalent ways of attaching propagators. We will now describe how to identify higher order terms in the vacuum functional by expanding in orders of $`\lambda `$ and $`\mathrm{}`$ with a recursive process using vacuum expectation values.
We begin with the lowest order expectation value, which is the tree level three field expectation value of order $`\lambda \mathrm{}^2`$. A covariant field theory calculation implies that this is
$$\begin{array}{cc}\hfill \mathrm{\hspace{0.17em}0}|\varphi (𝐱,0)\varphi (𝐲,0)\varphi (𝐳,0)|\mathrm{\hspace{0.17em}0}=i\lambda \mathrm{}^2\underset{\mathrm{}}{\overset{\mathrm{}}{}}dt\underset{\mathrm{}}{\overset{\mathrm{}}{}}& \underset{j=1}{\overset{3}{}}\mathrm{d}t_jW(t;t_1,t_2,t_3)\hfill \\ \hfill \times & G_𝐱(0;t_1)G_𝐲(0;t_2)G_𝐳(0;t_3)\hfill \end{array}$$
(2.7)
Alternatively, we can abandon the explicit kernel $`W`$ and think of this as some Feynman diagram which ties together three legs ending at $`x^0=0`$, in some way (the usual Feynman diagram showing three propagators meeting at a point vertex is not suitable):
(2.8)
Introducing an unknown diagram $`\mathrm{\Gamma }`$ into the vacuum,
$$\text{},$$
(2.9)
we require the expectation value to be given by
$$\mathrm{\hspace{0.17em}0}|\varphi (𝐱,0)\varphi (𝐲,0)\varphi (𝐳,0)|\mathrm{\hspace{0.17em}0}=𝒟\varphi \varphi (𝐱)\varphi (𝐲)\varphi (𝐳)\mathrm{\Psi }_0^2[\varphi ].$$
(2.10)
Equivalently this may be expressed as
(2.11)
where the fields in $`\mathrm{\Gamma }`$ have been contracted with the fields at external points using the inverse of the quadratic piece of the vacuum functional squared, which is the equal time propagator (the arcs on the right hand side). The factor of $`2`$ appears because of this squaring.
To identify $`\mathrm{\Gamma }`$ we invert these using (2.3). The right hand side of the above becomes the unknown kernel. The left hand side has external legs attached to the inverse of the equal time propagator. This produces propagators with a time derivative at the boundary $`t=0`$, and is independent of the sign of the time $`t_i`$ at the other end of the propagator (see for details). The result is
$$\text{}.$$
(2.12)
Note the modulus signs on the right hand side, illustrating the independence of the propagator endpoint on the sign of the time. In terms of the interaction kernel $`W`$ the vacuum state functional is
$$\mathrm{\Psi }_0[\varphi ]=\mathrm{\Psi }_0^{\text{free}}[\varphi ]\left(1\frac{\lambda }{2.3!\mathrm{}}\underset{\mathrm{}}{\overset{\mathrm{}}{}}dt\underset{j=1}{\overset{3}{}}\mathrm{d}t_jW(t;t_1,t_2,t_3)\underset{j=1}{\overset{3}{}}\varphi _𝐱G_𝐱(\stackrel{}{0};|t_j|)\right).$$
(2.13)
The tree level $`n`$–field term in the vacuum functional can be constructed similarly from the tree level $`n`$–field expectation value. Let us now describe a one loop term. We construct the vacuum functional tadpole using
$$\mathrm{\hspace{0.17em}0}|\varphi (𝐱,0)|\mathrm{\hspace{0.17em}0}=i\frac{\lambda \mathrm{}}{2!}dt\underset{j=1}{\overset{3}{}}\mathrm{d}t_jW(t;t_1,t_2,t_3)G_𝐱(0;t_1)G(t_2;t_3).$$
(2.14)
This is the generalisation of the tadpole
(2.15)
in local cubic theory. We illustrate it by
$$\text{}.$$
(2.16)
We have explicitly written in the factor of $`1/2!`$ normally implicit in a Feynman loop. The diagram is of order $`\mathrm{}\lambda `$, and receives contributions from the single field one loop term in the vacuum state and from the three field tree level piece we constructed above.
We include a new unknown in the vacuum,
$$\text{},$$
(2.17)
and calculating the expectation value implies
$$\text{}.$$
(2.18)
Inverting the equal time propagators as we did before identifies the new kernel,
$$\text{}.$$
(2.19)
Explicitly,
$$\begin{array}{cc}\hfill \mathrm{\Psi }_0[\varphi ]=\mathrm{\Psi }_0^{\text{free}}[\varphi ](1& +\frac{\lambda }{2.2!}dt\underset{j=1}{\overset{3}{}}\mathrm{d}t_jW(t;t_1,t_2,t_3)\varphi _𝐱G_𝐱(\stackrel{}{0};|t_1|)G(t_2;t_3)\hfill \\ \hfill +\frac{\lambda }{2.2!}dt\underset{j=1}{\overset{3}{}}\mathrm{d}t_j& W(t;t_1,t_2,t_3)\varphi _𝐱G_𝐱(\stackrel{}{0};|t_1|)G(\stackrel{}{0};|t_2|)G(0;0)G(\stackrel{}{0};|t_3|)\hfill \\ & \frac{\lambda }{2.3!\mathrm{}}\underset{\mathrm{}}{\overset{\mathrm{}}{}}\mathrm{d}t\underset{j=1}{\overset{3}{}}\mathrm{d}t_jW(t;t_1,t_2,t_3)\underset{k=1}{\overset{3}{}}\varphi _𝐱G_𝐱(\stackrel{}{0};|t_k|))\hfill \end{array}$$
(2.20)
This is very similar to the result found for the local theories in . Now we repeat the above arguments to construct the string field vacuum. The free field vacuum is (suppressing the $`\mathrm{}`$ dependence)
$$\text{},$$
(2.21)
just as in (2.4). The double represents the open or closed string field propagator. We propose the lowest order expansion
$$\text{}.$$
(2.22)
Trying to recover the three field vacuum expectation value using the vacuum functional as we did before, we find
(2.23)
The diagram on the left is the three field vacuum expectation value. As before we invert the equal time propagators on the right hand side of the above, and the left hand side becomes the kernel we are trying to identify. The first order cubic term in the vacuum state functional is therefore
(2.24)
Our gluing rules give us the result of joining propagators when each end lies on a constant time surface, they do not apply when one end lies on a time $`X^0(\sigma )`$ for non-trivial sigma dependence, nor how to attach more general surfaces together. For this we need Carlip’s method, which sheds some light on what the above kernel is.
In it was shown that to correctly sew the moduli spaces of two worldsheets the inverse propagator (first quantised Hamiltonian) should be attached to one of the boundaries being sewn. This removes a divergent integral over a redundant length parameter when we integrate over all shared field arguments on the boundaries being sewn.
Let $`T[t_1^{},t_2^{},t_3^{}]`$ is the three string amplitude with external legs ending at curves $`𝐗_i(\sigma )`$ at times $`X^0=t_i^{}`$ for $`i=1\mathrm{}3`$. Using the simple identity
$$f(0)=dsf(s)\delta (s)=dsdqf(s)G^1(s,q)G(q,0).$$
(2.25)
we can write the three field expectation value as
$$T[0,0,0]=\underset{j=1}{\overset{3}{}}\mathrm{d}t_j^{}\mathrm{d}t_jT[t_1^{},t_2^{},t_3^{}]\underset{i=1}{\overset{3}{}}G^1(t_i^{},t_i)G(t_i,0)$$
(2.26)
Now when we attach the inverse of the equal time propagator the gluing rules imply that the vacuum functional is
$$\mathrm{\Psi }_0[\mathrm{\Phi }]=\mathrm{\Psi }_0^{\text{free}}[\mathrm{\Phi }]\left(1\frac{i}{2.3!}\underset{j=1}{\overset{3}{}}\mathrm{d}t_j^{}\mathrm{d}t_jT[t_1^{},t_2^{},t_3^{}]\underset{i=1}{\overset{3}{}}G^1(t_i^{},t_i)\mathrm{\Phi }_𝐗G_𝐗(\stackrel{}{0};|t_k|)\right).$$
(2.27)
In our case the inverse propagator is that part of the first quantised string theory Hamiltonian which depends on $`𝐗`$ and $`t`$ on the boundary of our worldsheet. We can represent the vacuum functional as
(2.28)
where we have written $`H`$ in place of the inverse propagators for clarity. As before, the first order one loop term in the vacuum wave functional follows, and is
$$\text{}.$$
(2.29)
This construction applies to both open and closed strings.
## 3 The Schrödinger representation
In the Schrödinger representation $`n`$–field correlation functions are expressed by $`n`$ functional integrals,
$$\begin{array}{cc}\hfill \mathrm{\hspace{0.17em}0}|\varphi (𝐱_n,t_n)\mathrm{}\varphi (𝐱_1,t_1)|\mathrm{\hspace{0.17em}0}=𝒟& (\varphi _n\mathrm{}\varphi _1)\mathrm{\Psi }_0[\varphi _n]\varphi _n(𝐱_n)𝒮[\varphi _n,\varphi _{n1};t_nt_{n1}]\varphi _{n1}\hfill \\ & \mathrm{}\varphi _2(𝐱_2)𝒮[\varphi _2,\varphi _1;t_2t_1]\varphi _1(𝐱_1)\mathrm{\Psi }_0[\varphi _1],\hfill \end{array}$$
(3.1)
where there is one instance of the Schrödinger functional $`𝒮`$ for each product of two fields, and two instances of the vacuum state wave functional $`\mathrm{\Psi }_0`$ at the initial and final times (though the time dependence is trivial since the vacuum is an eigenstate. The Schrödinger functional is defined by
$$𝒮[\varphi _2,\varphi _1;t_2t_1]=\varphi _2|\mathrm{exp}\left(\frac{i}{\mathrm{}}\underset{t_1}{\overset{t_2}{}}dx^0H(x^0)\right)|\varphi _1.$$
(3.2)
The Feynman description is
$$𝒮[\varphi _2,\varphi _1;t]=𝒟\phi \mathrm{exp}\left(\frac{i}{\mathrm{}}\underset{0}{\overset{t}{}}dx^0L(\phi ,\dot{\phi })\right)|_{\phi =\varphi _1\text{ at }x^0=0,}^{\phi =\varphi _2\text{ at }x^0=t}$$
(3.3)
As before a shift of integration variable moves the field dependence out of the boundary conditions and into the action, giving
$$\begin{array}{cc}\hfill 𝒮[\varphi _2,\varphi _1;t]=𝒟\phi \mathrm{exp}(\frac{i}{\mathrm{}}\underset{0}{\overset{t}{}}\mathrm{d}x^0L(\phi ,\dot{\phi })& +\frac{i}{\mathrm{}}\mathrm{d}^D𝐱\dot{\phi }(𝐱,t)\varphi _2(𝐱)\hfill \\ & \frac{i}{\mathrm{}}\mathrm{d}^D𝐱\dot{\phi }(𝐱,0)\varphi _1(𝐱)\left)\right|_{\phi =0\text{ at }x^0=0}^{\phi =0\text{ at }x^0=t}\hfill \end{array}$$
(3.4)
The logarithm of the Schrödinger functional is the sum of connected Feynman diagrams constructed from a propagator $`G_D`$ which obeys Dirichlet boundary conditions at times $`x^0=0`$ and $`x^0=t`$ and with vertices integrated over the interval $`x^0[0,t]`$. All external legs end on the boundaries with a time derivative. The free field Schrödinger and vacuum functionals are therefore
$$\text{},$$
(3.5)
$$\text{},$$
(3.6)
where a propagator drawn with a dotted line obeys Dirichlet conditions on all boundaries (heavy lines) shown in the diagram. The gray dot denotes a time derivative, and an unbroken line will represent the free space propagator. The method of images allows us to write the propagators in these diagrams as sums over free space propagators. For example, the propagator in the vacuum is
$$G_d(𝐱_2,t_2;𝐱_1,t_1)=G_0(𝐱_2,t_2;𝐱_1,t_1)G_0(𝐱_2,t_2;𝐱_1,t_1),$$
in terms of the free space propagator $`G_0`$. Using this it is a simple matter to show the equality of (3.6) and (2.4).
The equivalence of Schrödinger and Heisenberg representations of quantum theories is well documented, but less well known are the details of this equivalence. In this section we use our diagrammatic methods to demonstrate how covariant correlation functions are given by the non-covariant Schrödinger representation. We return to the standard local $`\varphi ^3`$ interaction and calculate the simplest non-trivial correlation function
$$\text{}.$$
(3.7)
In the Schrödinger representation this is given by the functional integral
$$𝒟(\phi _2,\phi _1)\mathrm{\Psi }_0[\phi _2]\phi _2(𝐱_3)\phi _2(𝐱_2)𝒮[\phi _2,\phi _1;t]\phi _1(𝐱_1)\mathrm{\Psi }_0[\phi _1].$$
(3.8)
We need the first order Schrödinger and vacuum functionals, which are
(3.9)
$$\text{}.$$
(3.10)
The gray letters will be used to reference the diagrams. Excluding the loop diagrams, which can only contribute disconnected graphs, we must include all of the diagrams in (3.9) and (3.10) to calculate (3.7), not just diagram B which pictorially matches the desired result<sup>1</sup><sup>1</sup>1The reason is that, besides being of the correct order to contribute, the boundary conditions on the Dirichlet propagators mean that even though they may not end on one of the boundaries they can still see it, and can contribute to the free space result.
Inserting expressions (3.9) and (3.10) into (3.8) and carrying out the functional integrals the terms A to E contribute the set of diagrams shown below (excluding loops or disconnected pieces generated by the integrations). To illustrate we will give the calculation for diagram C explicitly, and state the results for the remaining diagrams. The numbers $`1`$ to $`3`$ will indicate the spatial positions of the external legs. The thick gray line is shorthand for $`K^1`$, the inverse of the quadratic piece of the product of the Schrödinger functional and vacuum wave functional which contracts indices in the first functional integral of (3.8),
We now begin the computation of diagram C. From (3.8) the $`\phi _1`$ integral does not see the contribution from the interacting Schrödinger functional, and is Gaussian with the insertion $`\phi _1(𝐱_1)`$,
The $`\phi _2`$ integral becomes
To obtain a connected graph we must contract $`\phi _2`$, $`\phi _3`$ with the 3 point function and the remaining 3 point function field must be contracted with the field attached to $`𝐱_1`$. There are $`3!`$ equivalent ways of doing this, so the result is
$$\text{}.$$
Using the method of images the propagator $`G_D`$ attached to one boundary is
(3.11)
From which the gluing properties tell us that the two types of terms we encounter in diagram C are
(3.12)
and
$$\text{}.$$
(3.13)
Defining the following sets of diagrams,
$$\text{},$$
(3.14)
$$\text{}.$$
(3.15)
Diagram C can be written
$$i\lambda \mathrm{}^2\underset{0}{\overset{t}{}}dz\mathrm{d}^D𝐳F_1(𝐳,z)H_2(𝐳,z)H_3(𝐳,z).$$
The calculation of the remaining diagrams A to E are similar and the results are below. There is no content in the length of the $`K^1`$ lines which varies in the diagrams only for clarity. $`K^1`$ will appear on the lower boundary and the equal time propagator on the upper boundary if we perform the $`\phi _1`$ integral first, and vice versa if we perform the $`\phi _2`$ integral first.
Diagram A contributes
Written in terms of the sums $`F_i(𝐳,z)`$ and $`H_i(𝐳,z)`$ this is (the order of the terms is respective to that of the diagrams above)
$`i\lambda \mathrm{}^2{\displaystyle \underset{0}{\overset{t}{}}}dz`$ $`{\displaystyle \mathrm{d}^D𝐳\left(G_0(𝐱_1,0;𝐳,z)G_0(𝐱_1,2t;𝐳,z)\right)H_2(𝐳,z)\left(H_3(𝐳,z)+G_0(𝐳,t;𝐱_3,0)\right)}`$
$`+(23)`$
$`+F_1(𝐳,z)\left(H_2(𝐳,z)+G_0(𝐳,z;𝐱_2,t)\right)\left(H_3(𝐳,z)+G_0(𝐳,z;𝐱_3,t)\right)`$
$`+(F_1(𝐳,z)+G_0(𝐱_1,2t;𝐳,z))(H_2(𝐳,z)+G_0(𝐳,z;𝐱_2,t))H_3(𝐳,z)+(23).`$
Diagram B contributes
which can be written
$`i\lambda \mathrm{}^2{\displaystyle \underset{0}{\overset{t}{}}}dz`$ $`{\displaystyle \mathrm{d}^D𝐳\left(G_0(𝐱_1,0;𝐳,z)G_0(𝐱_1,2t;𝐳,z)\right)H_2(𝐳,z)H_3(𝐳,z)}`$
$`+\left(F_1(𝐳,z)+G_0(𝐱_1,2t;𝐳,z)\right)H_2(𝐳,z)H_3(𝐳,z)`$
$`+F_1(𝐳,z)(H_2(𝐳,z)+G_0(𝐳,z;𝐱_2,t)(H_2(𝐳,z)+G_0(𝐳,z;𝐱_2,t)H_3(𝐳,z)`$
$`+(23).`$
Diagram D gives
which is
$$\begin{array}{cc}\hfill i\lambda \mathrm{}^2\underset{0}{\overset{t}{}}dz\mathrm{d}^D𝐳\left(F_1(𝐳,z)+G_0(𝐱_1,0;𝐳,z)\right)& \left(H_2(𝐳,z)+G_0(𝐳,z;𝐱_2,t)\right)\hfill \\ & \times \left(H_3(𝐳,z)+G_0(𝐳,z;𝐱_3,t)\right).\hfill \end{array}$$
Before calculating the diagrams coming from the interacting vacuum functional it is worthwhile adding the above terms up to find, with no calculation other than cancellations, that the total contribution from the interacting Schrödinger functional and the free vacuum is
$$i\lambda \mathrm{}^2\underset{0}{\overset{t}{}}dz\mathrm{d}^D𝐳G_0(𝐱_1,0;𝐳,z)G_0(𝐳,z;𝐱_2,t)G_0(𝐳,z;𝐱_3,t).$$
(3.16)
This is the correlation function we were looking for (with the correct coefficient) but the vertex is integrated in time only over the interval $`[0,t]`$. What remains must come from the vacuum. The contribution from the interacting vacuum functional $`\mathrm{\Psi }_0[\phi _1]`$, diagram E, is
$$\text{}.$$
The upper boundary (at time $`t`$) is broken as a reminder that only the outermost Dirichlet propagators vanish on both boundaries (i.e. come from the free field Schrödinger functional), whereas the inner three propagators vanish only at $`x^0=0`$ (come from the free field vacuum). This is
$$i\lambda \mathrm{}^2\underset{\mathrm{}}{\overset{0}{}}dz\mathrm{d}^D𝐳G_0(𝐱_1,0;𝐳,z)G_0(𝐳,z;𝐱_2,t)G_0(𝐳,z;𝐱_3,t),$$
(3.17)
There is a contribution from the vacuum $`\mathrm{\Psi }_0[\phi _2]`$, which we call term F,
Again, only the rightmost propagator comes from the Schrödinger functional. This is
$$i\lambda \mathrm{}^2\underset{\mathrm{}}{\overset{0}{}}dz\mathrm{d}^D𝐳G_0(𝐱_1,z;𝐳,t)G_0(𝐳,z;𝐱_2,0)G_0(𝐳,z;𝐱_3,0).$$
(3.18)
Changing variable
$$zz+t$$
(3.19)
in $`F`$ we have a total contribution of
$$\begin{array}{cc}\hfill i\lambda \mathrm{}^2\underset{\mathrm{}}{\overset{0}{}}dz\mathrm{d}^D𝐳& G_0(𝐱_1,0;𝐳,z)G_0(𝐳,z;𝐱_2,t)G_0(𝐳,z;𝐱_3,t)\hfill \\ & i\lambda \mathrm{}^2\underset{t}{\overset{\mathrm{}}{}}dz\mathrm{d}^D𝐳G_0(𝐱_1,0;𝐳,z)G_0(𝐳,z;𝐱_2,t)G_0(𝐳,z;𝐱_3,t)\hfill \end{array}$$
(3.20)
This is again the integrand we want, but we are missing the integral over the region $`[0,t]`$. This missing piece is precisely what we found to be contributed by the Schrödinger functional in (3.16).
We have seen that for this simple scattering process the interacting vacuum wave functional gives almost the correct result, up to some ‘small’ missing piece (in the sense that term is integrated over only a finite time), and that this is given by the Schrödinger functional.
### Conclusions
We have described, in analogy with quantum field theory, how to construct the string field theory vacuum wave functional for both open and closed strings. Note that we have not given, nor do we need to give, the explicit form of the three string vertex. That choice is arbitrary, the required kernels could even, theoretically, be computed with the Polyakov integral on the relevant manifold. For example, for the open string kernel $`\mathrm{\Gamma }`$ introduced in (2.22) this would be a disk with marked sections on the boundary. Although we have only given a few terms in the expansion of the vacuum functional thought we expect the construction to generalise to all orders in perturbation theory.
Given the gluing property for string fields we may postulate a string field theory interaction which is non-local in the spatial co-ordinates, but local in time, and where the fields depend only on $`t`$ and the set $`𝐗`$ defined in the introduction. For example we could define an interaction very like that in the light-cone gauge , with worldsheet $`\tau `$ replaced by spacetime $`t`$,
$$\begin{array}{cc}\hfill \underset{\mathrm{}}{\overset{\mathrm{}}{}}\mathrm{d}t𝒟(𝐗,𝐗_1,𝐗_2)\mathrm{\Phi }[& 𝐗(\sigma ),t]\mathrm{\Phi }[𝐗_1(\sigma _1),t]\mathrm{\Phi }[𝐗_2(\sigma _2),t]\hfill \\ & \underset{0\sigma \pi /2}{}\delta (𝐗(\sigma )𝐗_1(\sigma _1))\underset{\pi /2\sigma \pi }{}\delta (𝐗(\sigma )𝐗_2(\sigma _2))\hfill \end{array}$$
(3.21)
where all three fields live at the same time (as in a particle theory) and the parameters along the string are
$$\begin{array}{cc}\hfill \sigma _1& =2\sigma ,0\sigma \pi /2\hfill \\ \hfill \sigma _2& =2\sigma \pi ,\pi /2\sigma \pi \hfill \end{array}$$
(3.22)
so that each string has parameter domain $`0\sigma \pi `$. If the free Hamiltonian was the inverse of $`G(𝐗_2,t_2;𝐗_1,t_1)`$ we could repeat any quantum field theory argument in this theory. However, this approach removes the $`X^0`$ oscillators from the theory at the outset, and it is unclear how to justify this nor how amplitudes in this theory relate to known results.
The interaction is local in time, but as with the light-cone gauge we do not expect locality in one direction to spoil the UV finiteness of string interactions . Our time co-ordinate, unlike in other string field theories, is a time at which the whole spatially extended string (with ghosts) exists. This seems to work around the problems of, for example, quantising Witten’s theory where time is normally taken to be the midpoint of $`X^0(\sigma )`$, but the string remains extended in time. This may be a worthwhile avenue for future study.
Although it is a lengthier task, it is not especially more difficult to reconstruct the Schrödinger functional of quantum field theory using a similar approach, and the generalisation to string field theory will follow. We would propose the first order expansion
$$𝒮[\varphi _2,\varphi _1,t]=𝒮^{\text{free}}[\varphi _2,\varphi _1,t]\left(1+\varphi _2\varphi _2\varphi _2𝒮^{(3,0)}+\varphi _2\varphi _2𝒮^{(2,1)}\varphi _1+\varphi _2𝒮^{(1,2)}\varphi _1\varphi _1+𝒮^{(0,3)}\varphi _1\varphi _1\varphi _1\right)$$
(3.23)
which includes four unknown kernels, $`𝒮^{(i,j)}`$. There are four possible three field amplitudes which must be reproduced by a functional integration with one instance the Schrödinger functional,
$`\mathrm{\hspace{0.17em}0}|\varphi (𝐱,t)\varphi (𝐲,t)\varphi (𝐳,t)|\mathrm{\hspace{0.17em}0}`$ $`\mathrm{\hspace{0.17em}0}|\varphi (𝐱,t)\varphi (𝐲,t)\varphi (𝐳,0)|\mathrm{\hspace{0.17em}0}`$
$`\mathrm{\hspace{0.17em}0}|\varphi (𝐱,t)\varphi (𝐲,0)\varphi (𝐳,0)|\mathrm{\hspace{0.17em}0}`$ $`\mathrm{\hspace{0.17em}0}|\varphi (𝐱,0)\varphi (𝐲,0)\varphi (𝐳,0)|\mathrm{\hspace{0.17em}0}`$
giving us four simultaneous functional equations with which to determine the kernels $`𝒮^{(i,j)}`$.
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# RECENT RESULTS FROM KTEV
## 1 $`K_L^0\pi ^0\mu ^+\mu ^{}`$
The decay $`K_L^0\pi ^0\mu ^+\mu ^{}`$ has three contributions to the amplitudes: a $`CP`$ conserving term from intermediate states with two photons; a $`CP`$ violating term from indirect $`CP`$ violation; and a direct $`CP`$ violating term from 2nd-order electroweak penguin and box diagrams. The process continues to be of theoretical interest; see in particular the presentation by Christopher Smith at this conference. There is also interest on the experimental side as a consequence of a recent unusual result from the HyperCP collaboration $`^\mathrm{?}`$.
In searching for $`\mathrm{\Sigma }^+p\mu ^+\mu ^{}`$, HyperCP found three events. This is not so unusual; what is peculiar is that to within the $`0.5MeV`$ resolution of the mass measurements, all three events have the same dimuon mass of 214.3$`MeV`$. In the standard model, the muon pair is produced by an off-shell photon and a spread of dimuon masses is expected. The HyperCP collaboration estimates that the probability of three events from an intermediate photon having the same mass to this level of precision is about 0.8%, and they suggest that there may be a new intermediate neutral state causing this anomaly.
This postulated new state would be a flavor changing neutral current when coupling to quarks, but not when coupling to leptons. The scale of the partial width for the three observed events, $`\mathrm{\Gamma }(\mathrm{\Sigma }^+pP^0,P^0\mu ^+\mu ^{})2.6\times 10^{19}MeV`$ is too small for a strong interaction, which have widths on the order of a few $`10^{12}MeV`$ in $`\mathrm{\Sigma }^+`$ decays. Also of course, the active state of research into new bound QCD states not withstanding, a new narrow hadronic resonance in this mass range would be a surprise. So we consider new point-like particles. Allowing that the new interaction conserves parity and angular momentum, the only possible $`(J)^P`$ values for decays into $`\mu ^+\mu ^{}`$ are $`(0)^{}`$and $`(1)^{}`$. However if this $`P^0`$ is a vector current then it will also contribute to $`K_L^0\pi ^0\mu ^+\mu ^{}`$decay.
The existing KTeV limit $`^\mathrm{?}`$ on Br($`K_L^0\pi ^0\mu ^+\mu ^{}`$) of $`3.8\times 10^{10}`$ at the 90% C.L., while an order of magnitude above the standard model prediction $`^\mathrm{?}`$ of $`(1.5\pm 0.3)\times 10^{11}`$, does correspond to a partial width of $`\mathrm{\Gamma }(K_L^0\pi ^0P^0,P^0\mu ^+\mu ^{})4.8\times 10^{24}MeV`$, nearly 5 orders of magnitude less than the postulated HyperCP rate. The vector current hypothesis is thus disfavored.
Work on a new limit for Br($`K_L^0\pi ^0\mu ^+\mu ^{}`$) based on the full KTeV data set is in progress.
## 2 $`\mathrm{\Xi }^0\mathrm{\Sigma }^+\mu ^{}\nu `$
The first observation $`^\mathrm{?}`$ of the decay $`\mathrm{\Xi }^0\mathrm{\Sigma }^+\mu ^{}\nu `$ was made on the 1997 KTeV dataset; here we report results based on the 1999 data, which corresponds to about $`3\times 10^8`$ $`\mathrm{\Xi }^0`$ decays. We reconstruct the $`\mathrm{\Sigma }^+`$ hyperon in the $`p\pi ^0;\pi ^0\gamma \gamma `$ mode and normalize the data sample with $`\mathrm{\Xi }^0\mathrm{\Lambda }\pi ^0;\mathrm{\Lambda }p\pi ^{}`$. Backgrounds are from $`\mathrm{\Xi }^0`$, $`\mathrm{\Lambda }`$ and $`K_L^0\pi ^+\pi ^{}\pi ^0`$ decays; for all decays other than $`\mathrm{\Xi }^0\mathrm{\Lambda }\pi ^0;\mathrm{\Lambda }p\pi ^{}`$ and $`K_L^0\pi ^+\pi ^{}\pi ^0`$, Monte Carlo samples corresponding to 10 or more times the data sample have been generated. The kaon background is studied with data events where the highest-momentum track is negatively charged. For $`K_L^0\pi ^+\pi ^{}\pi ^0`$, high momentum $`\pi ^+`$ and $`\pi ^{}`$ are equally probable, but the hyperon signal is overwhelmingly comprised of events with high-momentum positively charged tracks. Background from $`\mathrm{\Xi }^0\mathrm{\Lambda }\pi ^0;\mathrm{\Lambda }p\pi ^{}`$ is suppressed by relying on the neutrino in the signal mode to produce a missing momentum component perpendicular to the line of $`\mathrm{\Xi }^0`$ flight. The overall background level is very low, and there are nine events in the data, as shown in Figure 1. We obtain, as a preliminary result, Br($`\mathrm{\Xi }^0\mathrm{\Sigma }^+\mu ^{}\nu `$) $`=(4.3\pm 1.4)\times 10^6`$.
## 3 $`|V_{us}|`$
A long standing issue in flavor physics has been a discrepancy, at the 2$`\sigma `$ level, of the measured values of $`|V_{ud}|`$, $`|V_{us}|`$ and $`|V_{ub}|`$ from the unitarity constraint for the first row of the CKM matrix. It should be noted that the first row of the matrix is the one that provides the most stringent test of 3-generation unitarity. The value of $`|V_{ud}|`$ is precisely known from nuclear and neutron beta decays, and $`|V_{ub}|`$ is small enough to be irrelevant here. $`|V_{us}|`$ may be determined from semileptonic kaon decays; in the past only the $`K_L^0\pi ^0e^\pm \nu `$ mode has been used, due to uncertainty in the form factors of the $`K_L^0\pi ^0\mu ^\pm \nu `$ mode.
To address this discrepancy, KTeV has sought improved measurements of $`\mathrm{\Gamma }`$($`K_L^0\pi ^0e^\pm \nu `$) and $`\mathrm{\Gamma }`$($`K_L^0\pi ^0\mu ^\pm \nu `$). There being no other decay modes that are known to sufficient precision to normalize the data set, we have measured five ratios of partial widths for the modes $`K_L^0\pi ^0e^\pm \nu `$, $`K_L^0\pi ^0\mu ^\pm \nu `$, $`K_L^0\pi ^+\pi ^{}\pi ^0`$, $`K_L^03\pi ^0`$, $`K_L^0\pi ^+\pi ^{}`$and $`K_L^02\pi ^0`$. These modes account for nearly all $`K_L^0`$ decays, and this fact may be used to extract the branching ratios for the semileptonic modes. With the previously accepted $`^\mathrm{?}`$ value (see the presentation of Gaia Lanfranchi at this conference) for the lifetime of the $`K_L^0`$, partial widths for the semileptonic modes and then values for $`|V_{us}|`$ have been extracted.
A full exposition of this work and the related analyses spans a number of publications:
* The radiative corrections - reference $`^\mathrm{?}`$
* Measurement of form factors for $`K_L^0\pi ^0e^\pm \nu `$ and $`K_L^0\pi ^0\mu ^\pm \nu `$ \- reference $`^\mathrm{?}`$
* Check with radiative $`K_L^0\pi ^0e^\pm \nu `$ and $`K_L^0\pi ^0\mu ^\pm \nu `$ decays - reference $`^\mathrm{?}`$
* The measurement of the partial width ratios - reference $`^\mathrm{?}`$
* The extraction of $`|V_{us}|`$ from these - reference $`^\mathrm{?}`$
The final result, using $`f_+(0)=0.961\pm 0.008`$, is $`|V_{us}|`$ $`=0.2252\pm 0.0008_{\mathrm{KTeV}}\pm 0.0021_{\mathrm{ext}}`$, where the $`\mathrm{KTeV}`$ uncertainty includes uncertainties in the KTeV branching fractions and form factor measurements, and the $`\mathrm{ext}`$ uncertainty includes uncertainties in $`f_+(0)`$, $`K_L^0`$ lifetime, and radiative corrections. This result does resolve the unitarity discrepancy, and the attendant results show a high degree of internal consistency. While our branching ratios are not in good agreement with the values listed in the 2002 PDG, many of the discrepancies may be explained by postulating that Br($`K_L^0\pi ^0e^\pm \nu `$) in the 2002 PDG was too low. In discussion I emphasized that, as noted by Cirigliano et.al. $`^\mathrm{?}`$, the PDG did not actually have a precise and direct measurement of Br($`K_L^0\pi ^0e^\pm \nu `$) on hand. The value that they recommended was inferred from reports of other measurements and the indisputable constraint that the sum of the branching ratios must be one. It was, if you would, a ”global fit”, albeit a substantially simpler analysis than many calculations going by that name. Caveat!
## 4 Acknowledgements
I very much want to thank the organizers of this most excellent conference, and in particular the gracious and highly competent Elizabeth Hautefeuille as well as the sagacious and justly renowned Jean Tran Thanh Van.
The KTeV collaboration gratefully acknowledges the support and effort of the Fermilab staff and the technical staffs of the participating institutions for their vital contributions. This work was supported in part by the U.S. Dept. of Energy, the National Science Foundation, and the Ministry of Education and Science of Japan.
## References
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# Analysis of radiation-pressure induced mechanical oscillation of an optical microcavity
## Abstract
The theoretical work of V.B. Braginsky predicted that radiation pressure can couple the mechanical, mirror-eigenmodes of a Fabry-Pérot resonator to it’s optical modes, leading to a parametric oscillation instability. This regime is characterized by regenerative mechanical oscillation of the mechanical mirror eigenmodes. We have recently observed the excitation of mechanical modes in an ultra-high-Q optical microcavity. Here, we present a detailed experimental analysis of this effect and demonstrate that radiation pressure is the excitation mechanism of the observed mechanical oscillations.
The work of V.B. BraginskyBraginsky et al. (2001, 2002) predicted that pressure induced by circulating radiation in a Fabry-Pérot resonator can couple the optical modes to the mechanical mirror eigenmodes. The coupling can lead to a parametric oscillation instability, characterized by regenerative oscillation of the mechanical mirror eigen-modes. This mechanism has been studied theoretically for its possible role in setting a detection sensitivity limit in the laser interferometer gravitational wave observatory (LIGO) Abramovici et al. (1992); Amelino-Camelia (1999), but has so far not been observed experimentally.
Recently, we have observed a nonlinear mechanismRokhsari et al. (2004) in ultra-high-Q toroid microcavitiesArmani et al. (2003) that is distinct from other nonlinear mechanisms already observed in these structuresKippenberg et al. (2004); Spillane et al. (2002). The geometry for observation of this nonlinearity is a standard one in which a wave (here referred to as the pump) is coupled from a waveguide to a microcavity mode. The nonlinearity manifests itself as oscillations in the pump power transmitted past the micro-cavity. These oscillations are observed to occur at a distinct threshold pump power level and have spectral components at characteristic frequencies. Numerical modeling and spectral analysis reported in Ref. Rokhsari et al. (2004) revealed that the observed oscillations are due to regenerative oscillation of certain mechanical eigenmodes of the toroid microcavity. In this letter, we demonstrate that the observed mechanical oscillations are caused by radiation pressure, and specifically rule out another mechanism (thermal effectsZalalutdinov et al. (2001) ). As such, this work confirms the first observation of radiation-pressure-induced parametric oscillation instability as predicted by Braginsky.
The theoretical treatment of Braginsky Braginsky et al. (2001) considered mechanical oscillations of Fabry-Pérot mirror eigenmodes which leads to Stokes and anti-Stokes sidebands (at frequencies, $`\omega _0\pm \omega _m`$ where $`\omega _0`$ is the optical and $`\omega _m`$ the mechanical frequency). It was shownBraginsky et al. (2001) that if the Stokes field coincides with an adjacent optical cavity mode the phenomenon of parametric oscillation instability can occur. In contrast to the Braginsky theory, we observed mechanical oscillations of several mechanical modes (above a certain threshold) when the mechanical resonance frequencies ($`\omega _m`$) produce Stokes and anti-Stokes fields that fall within the same cavity resonance (i.e. $`\omega _m<\frac{\omega _0}{Q}`$)Rokhsari et al. (2004). For Q-factors in the range of 10<sup>6</sup>-10<sup>8</sup> this corresponds to frequencies in the range of ca. 1-100 MHz which coincides with the range of the first three fundamental mechanical modes of the toroid microcavities employed in this work. Fig. 1a shows the first, three mechanical modes of a toroid microcavity and Fig. 1b their frequency dependence on cavity length. Note that the mechanical motion causes modulation of the optical pathlength of the toroid cavity modes, causing the excitation of optical sidebands. These fields appear in the cavity transmission spectrum, as shown in Fig.1c.
To account for this scenario we have extended the coupled mode analysis of Braginsky to the present case of Stokes and anti-Stokes frequency pairs falling within the same cavity resonance (For simplicity only one pair is considered here). In addition optical coupling effects associated with the waveguide-resonator junction are, by necessity, included in the analysis. Using the rotating wave and the slowly varying envelope approximation for all field amplitudes, the mutual coupling of the pump $`(a_p),`$ Stokes $`(a_S)`$ anti-Stokes, $`(a_{AS})`$ and mechanical mode $`(x_m)`$, can be described by the following coupled-mode equations:
$`{\displaystyle \frac{x_m}{t}}`$ $`={\displaystyle \frac{1}{2\tau _m}}x_m+{\displaystyle \frac{iK_{om}}{2\sqrt{m_{eff}}C(\mathrm{\Gamma })}}(a_p^{}a_{AS}+a_pa_S^{})`$ (1)
$`{\displaystyle \frac{a_p}{t}}`$ $`={\displaystyle \frac{a_p}{2\tau }}+i\mathrm{\Delta }\omega a_p+{\displaystyle \frac{iK_{mo}}{\sqrt{m_{eff}}\omega _m}}(x_m^{}a_{AS}+x_ma_S)+\kappa s`$
$`{\displaystyle \frac{a_S^{}}{t}}`$ $`={\displaystyle \frac{1}{2\tau }}a_S^{}+i\left(\mathrm{\Delta }\omega \omega _m\right)a_S^{}{\displaystyle \frac{iK_{mo}}{\sqrt{m_{eff}}\omega _m}}x_ma_p^{}`$
$`{\displaystyle \frac{a_{AS}}{t}}`$ $`={\displaystyle \frac{1}{2\tau }}a_{AS}+i\left(\mathrm{\Delta }\omega +\omega _m\right)a_{AS}+{\displaystyle \frac{iK_{mo}}{\sqrt{m_{eff}}\omega _m}}x_ma_p`$
In these equations, the optical pump is detuned from the cavity-mode line-center by $`\mathrm{\Delta }\omega =\omega \omega _0`$. The Stokes and anti-Stokes frequencies lie within the resonance bandwidth of the pump mode, and, correspondingly, are detuned by $`\mathrm{\Delta }\omega _{AS}=\left(\mathrm{\Delta }\omega +\omega _m\right)`$ and $`\mathrm{\Delta }\omega _S=\left(\mathrm{\Delta }\omega \omega _m\right)`$ . The first equation describes the mechanical eigenmode where $`\left|x_m\right|^2`$ is normalized to mechanical energy, i.e. $`E=_iϵ_i\sigma _i𝑑V`$, (where $`ϵ_i`$ and $`\sigma _i`$ are the diagonal components of the strain and stress tensor) which decays with the lifetime $`\tau _m(Q_m=\omega _m\tau _m)`$. Correspondingly, $`\left|a_p\right|^2`$ is the energy in the pump mode, $`|s|^2`$ is the launched pump power in the waveguide. The total lifetime of the optical modes is given by $`\frac{1}{\tau }=\frac{1}{\tau _0}+\frac{1}{\tau _{ex}},`$where the external lifetime $`(\tau _{ex})`$ describes coupling of the microcavity mode to the waveguide via $`\kappa =i\sqrt{\frac{1}{\tau _{ex}}}`$ and $`K\frac{\tau _0}{\tau _{ex}}`$ is the normalized coupling constant. $`C`$($`\mathrm{\Gamma })`$ is a correction factor \[$`\mathrm{1..2}`$\] due to reduction of circulating power in the presence of modal couplingKippenberg et al. (2002)$`K_{mo}\frac{\omega _0}{R}`$ describes the mechanically-induced displacement of the optical cavity resonant frequency and contains, in general, a contribution from direct spatial change as well as refractive index changes (stress-optical effect) Str Shelby et al. (1985). The effective coupling of optical mode to the mechanical mode is governed by $`K_{om}\frac{1}{Rn_{eff}}`$ in the case of radiation pressureBraginsky et al. (2001). The effective mass $`m_{eff}`$ appearing in Eqn. (1) is calculated numerically, by evaluating the total mechanical energy $`E_m`$ in the mechanical mode and the corresponding harmonic radial displacement (amplitude $`r`$) of the toroid periphery wherein the optical mode circulates (compare fig. 4b)Mec . Solving the coupled mode equations in steady-state the pump-power threshold for onset of mechanical oscillations is given by:
$`P_{thresh}`$ $`=\left({\displaystyle \frac{\omega _m}{\omega _0}}\right)R^2m_{eff}{\displaystyle \frac{1}{\tau _m}}{\displaystyle \frac{1}{\tau _0}}{\displaystyle \frac{|1+K+2i\mathrm{\Delta }\omega \tau _0|^2}{4K}}`$ (2)
$`{\displaystyle \frac{1}{2\tau }}[{\displaystyle \frac{1}{1+4\tau ^2\mathrm{\Delta }\omega _{AS}^2}}{\displaystyle \frac{1}{1+4\tau ^2\mathrm{\Delta }\omega _S^2}}]^1`$
Careful inspection of the last term of the threshold equation shows that mechanical gain is possible (i.e., positive threshold power) for $`\omega >\omega _0(`$i.e. $`\mathrm{\Delta }\omega >0)`$. For $`\mathrm{\Delta }\omega <0`$, the mechanical mode is damped. The need to overcome mechanical loss leads to the $`\frac{1}{\tau _m}`$dependence, while the dependence of radiation pressure upon circulating optical power leads to the $`\frac{1}{\tau _0}`$-dependence as well as the presence of a weighting factor describing the effect of waveguide coupling $`K\frac{\tau _{ex}}{\tau _0}`$ and pump detuning $`\mathrm{\Delta }\omega `$. The optical-Q scaling dependences fall into two regimes. The first occurs when $`\omega _m<\frac{1}{\tau }`$ . In this regime the mechanical oscillation threshold exhibits an inverse cubic dependence on optical Q ($`P\frac{1}{Q_m}\left(\frac{1}{Q_0}\right)^3`$) . In contrast, for $`\omega _m>\frac{1}{\tau }`$ (herein called the high-frequency (HF) regime), the rapid $`1/Q_0^3`$ dependence is reduced because the Stokes field build-up is less-and-less effective in creating radiation pressure. In this regime, minimum threshold can be shown to occur over-coupled (i.e., $`K>1`$), where again the condition $`\omega _m<\frac{1}{\tau }`$ is met (i.e. the mechanical oscillation frequency is again less than the ”loaded” cavity bandwidth), which causes the minimum threshold (i.e., $`\frac{^2P}{K\mathrm{\Delta }\omega }=0`$ thr ) to approach an asymptotic value. The transition to the HF regime, under conditions of optimum threshold, occurs for an optical Q given by $`Q_0^{HF}\frac{\omega _0}{\omega _m}`$.
To confirm these theoretical predictions the threshold dependence (as given by Eq. 2) on both optical and mechanical Q-factor have been measured. The data presented are taken using a single microtoroid device. Coupling to the resonator was performed using a fiber-optic taper coupler (see inset fig.2). The micro-toroid under consideration had principal, pillar and minor diameters of $`72,36`$ and $`6.8\mu m`$, respectively, and possessed mechanical resonances frequencies at $`4.4`$, $`25.6`$ and $`49.8`$ MHz for the first three mechanical modes ($`n=1,2,3`$ and $`m=0`$). The optical pump wavelength was $``$1550 nm and mechanical oscillation instability was observed by detecting the characteristic oscillations in the transmitted pump power (compare Fig. 1c) Rokhsari et al. (2004) using an electrical spectrum analyzer (ESA) as described in Ref Rokhsari et al. (2004). Optimization of coupling ($`K`$) was performed by adjustment of the gap between the fiber taper and the microtoroid as described in refs. Spillane et al. (2002); Kippenberg et al. (2004). To measure the dependence of the oscillation threshold on $`Q_m`$, a silica microprobe was brought into contact with the interior (disk region) of the toroid structure. Variation of the probe contact force thereby modified mechanical Q while leaving the optical Q unaffected. The microprobe, which was made from an optical fiber, had a tip diameter of $``$2 $`\mu m`$ and can be seen in the inset of Fig. 2. In the absence of probe contact $`Q_m`$ was measured to be $`5000`$ for the $`n=1`$ mode, and upon progressive increase in tip pressure could be continuously decreased to $`Q_m`$ $`50`$. Below threshold, the thermal displacement of the mechanical eigenmodes (the temperature being $`300`$K) provides sufficient modulation to be optically detectable, causing the appearance of Lorenzian peaks in the cavity transmission spectrum. $`Q_m`$ was then determined by fitting the transmission spectrum with a Lorenzian, as shown in the inset of Fig. 2. For each $`Q_m`$, the minimum threshold was measured for the $`n=1`$ flexural mode as shown in Fig. 2. The solid line in the main panel shows that the data exhibit the $`1/Q_m`$ dependence in agreement with Eqn. 2, and Ref. Braginsky et al. (2001).
We next measured the threshold dependence on the optical Q factor as shown in Fig. 3 for both the $`n=1`$ (main panel) and the $`n=3`$ (inset) mechanical modes. The optical Q factor was adjusted by exciting different radial and transverse optical modes. For lower optical Q, wherein the acoustical oscillation frequency falls within the cavity bandwidth, the rapid $`1/Q^3`$ dependence is observed for $`n=1`$ as predicted. For higher optical Q, as theoretically predicted a transition into the HF regime occurs at $`Q_0^{HF}10^7`$. This point agrees well with the theoretical prediction ($`\frac{\omega _0}{\omega _m}`$). It is important to note that these observations rule out thermal effects Zalalutdinov et al. (2001) as origin of the observed oscillationsthe . In Fig. 3, the solid line is the minimum threshold i.e. ($`\frac{^2P}{K\mathrm{\Delta }\omega }=0`$) as given by equation (2). With the exception of the effective mass, $`m_{eff}`$, all parameters used to create this plot were experimentally measured parameters (i.e., $`C(\mathrm{\Gamma }),R,Q_m,\omega _m,Q_0,\omega _0)`$. The effective mass $`m_{eff}`$ was inferred to be $`m_{eff}^{(1)}`$ $`=3.3\times 10^8`$ $`kg`$.
The inset of Fig. 3 shows the measured threshold versus Q for the $`n=3`$ mode. The $`n=3`$ mode threshold dependence shows that this mode is already well into the HF regime, exhibiting the theoretically predicted asymptotic behavior of the minimum threshold. This fact is consistent with the observed resonance frequency, 49 MHz, for the $`n=3`$ mode which predicts that the HF regime occurs for optical Q factors in excess of 10<sup>7</sup> ($`Q_0^{HF}=\frac{\omega _0}{\omega _m}=3.8\times 10^6`$). Comparison with the $`n=1`$ mode data shows that oscillation on the $`n=3`$ mode is preferred for lower optical Qs. Indeed, preference to the $`n=3`$ mode was possible by loading the microcavity into the over-coupled regime, in agreement with theory. The solid curve in the inset gives the single-parameter fit to the $`n=3`$ data which yields $`m_{eff}^{(3)}=5\times 10^{11}kg`$, which is a factor of 660 lower than the mass of the $`n=1`$ mode.
As a further test of the validity of the theoretical model, the experimental effective mass values are compared with the theoretical prediction based on finite element modeling. For the $`n=3`$ mode, the predicted mass associated with the radial motion was $`m_{eff}^{(3)}=5\times 10^{11}kg`$, which is in very good agreement with the experimental fit. However, for the $`n=1,2`$ modes, the calculated effective mass is a strong function of the offset of the toroidal ring with respect to the equatorial plane of the diskDou . To both validate and quantify this offset, a cross section of the toroid microcavity used in this study was obtained with focused ion beam slicing. SEM imaging (cf. Fig. 4 panel A) reveals the presence of the above-postulated equatorial offset which amounts to $`\mathrm{\Delta }=1.3`$ $`\mu m`$ . Incorporation of this offset into the numerical mass calculation yields $`m_{eff}^{(1)}=2.6\times 10^8kg`$ and $`m_{eff}^{(2)}=2\times 10^9kg`$. This value agrees very well with the experimental values from above. Finally, the numerical model also explains why the $`n=2`$ mode is only observed sub-threshold in the experiments. The low mechanical Q value ( $`200`$) in conjunction with its high effective mass and frequency, predicts threshold powers $`>2`$ mW, which are higher than pump powers available in the experiments.
In summary, presented is both an experimental and theoretical analysis of radiation pressure induced parametric oscillation instability, as predicted by Braginsky. Excellent agreement of the threshold functional dependence on optical Q is obtained, providing a confirmation that radiation pressure is the excitation mechanism of the observed oscillations. Besides the fundamental aspects of this work, the observed coupling of mechanical and optical modes by radiation pressure can find applications in micro- and nanomechanical systems (MEMS/NEMS)Craighead (2000) for ultra-high sensitivity measurements of charge Cleland and Roukes (1998), displacement Rugar et al. (2004), mass, forceRugar et al. (2004) or biological entities Ilic et al. (2004). Equally important, radiation pressure as observed here can be used to achieve cooling of a mechanical resonator mode.
This work was supported by the NSF, DARPA and the Caltech Lee Center. T.J.K. acknowledges an IST-CPI postdoctoral fellowship.
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# Casimir Energies for 6D Supergravities Compactified on 𝒯₂/𝑍_𝑁 with Wilson Lines
## 1 Introduction.
The Casimir energy for various field theories compactified on 2 internal dimensions has been extensively studied (see for example refs. -). In this paper we return to this calculation for the particularly simple examples of 2-tori and their orbifolds. At the technical level, our aim in is to provide a generalization of previous calculations for fields compactified on $`𝒯_2`$ and $`𝒯_2/Z_N`$ to include very general boundary conditions, including those which would be generated by (constant) Wilson-line backgrounds.
Our physical motivation for performing this calculation is due to its relevance for the recent proposal for using 6D supergravity to shed light on the cosmological-constant problem within the context of Supersymmetric Large Extra Dimensions (SLED). (See also for related discussions and similar proposals.) Within this proposal the extra dimensions must presently be sub-millimeter in size, and the recently-discovered cosmological Dark Energy density corresponds to the Casimir energy of the model’s bulk fields as functions of the moduli of these large 2 internal dimensions. The smallness of the Dark Energy density in this picture is ultimately traced to the large size of these extra dimensions. In this picture these moduli can be extremely light in a technically natural way, and so the Dark Energy phenomenology is controlled by the dynamics of the moduli as they respond to this Casimir energy. Although simple cosmologies appear to be possible along roughly these lines , a more detailed determination of their viability requires a more careful calculation of the modulus potential which is provided by the Casimir energy, such as we present here.
Torus-based Casimir energies are particularly well-studied, and recent one-loop studies include , who study the energy of massless scalars compactified on $`𝒯_2`$ and $`𝒯_2/Z_2`$ at the special modular point corresponding to an orthogonal underlying torus. Ref. , on the other hand, computes the full modulus-dependence of the Casimir energy for various massless fields compactified on $`𝒯_2`$, assuming these fields to satisfy periodic boundary conditions about the cycles of the background geometry. Ref. computes for $`𝒯_2`$ compactifications the dependence of the Casimir energy on a particular modulus, arg$`U`$, but restrict some others (by choosing equal toroidal radii). Refs. and consider the $`𝒯_2/Z_2`$ orbifold, with moduli fixed to those values appropriate for an orthogonal underlying torus. Ref. considers the case of $`𝒯_2/Z_2`$ or $`𝒯_2`$, including the presence of Wilson lines, but only computes the Wilson-line dependent part of the result. Other recent calculations make similar assumptions .
In this paper we present results for the fields which appear in supergravity models compactified on a 2-torus, $`𝒯_2`$ and some of its orbifolds, $`𝒯_2/Z_N`$, as functions of the relevant moduli and the higher-dimensional particle mass, $`m`$. We do so for a very general set of boundary conditions about the cycles of the background geometry, such as could be generated by the presence of nontrivial Wilson lines wrapping these cycles. We provide formulae which lend themselves to numerical evaluation, and focus in particular on the form of the result in the large- and small-$`m`$ limits.
A spin-off of this calculation is the information it provides about the large-$`m`$ limit, and so of the ultraviolet sensitivity of these Casimir-energy calculations. This UV-sensitivity is a crucial part of the SLED proposal, and explicit calculations such as those presented here provide important checks on the more general, heat-kernel, UV-sensitivity calculations for 6D backgrounds, presented in a companion paper, . These general heat-kernel results properly describe the explicit dependence of the toroidal example considered here, and also show how flat geometries like tori are dangerous gedanken laboratories, because they are particularly insensitive to large-$`m`$ UV effects.
The explicit orbifold calculation also provides the simplest example of the appearance of new, brane-localized, UV divergences — a phenomenon which is generic to Casimir-energy calculations in the presence of branes due to the singularities and boundaries which these branes typically induce in the background geometry . We exhibit this new divergent contribution explicitly and show that it arises due to the orbifold projections which must be performed.
So far as they go, our results support the SLED framework inasmuch as all of the one-loop contributions we find are at most of order $`1/𝒜^2`$, where $`𝒜`$ is the compactification area. This makes them no larger than is required for the success of the SLED proposal.
The paper is organized as follows: In the next section we describe the relevant parts of the toroidal geometry and the boundary conditions for fields on this geometry for which we compute the Casimir energy. The explicit calculations are first performed in detail for complex scalar fields, for which we also display the ultraviolet-divergent and large-$`m`$ dependent parts. Results for the finite part of the Casimir energy are given for general boundary conditions for massless scalar fields and these are then extended using a simple argument to obtain the corresponding results for massless higher-spin fields in 6 dimensions. Section 3 addresses the same issues for the case of $`𝒯_2/Z_N`$ orbifolds, with the $`𝒯_2/Z_2`$, $`𝒯_2/Z_4`$ orbifolds treated in detail. The small-$`m`$/large-$`m`$ cases are again discussed. The main text only quotes the results of the calculations, and full technical details are provided in the Appendix. Many of the tools of this Appendix including the detailed calculations of series of Kaluza-Klein integrals, can prove useful for other applications such as loop corrections to gauge couplings in gauge theories on orbifolds with discrete Wilson lines.
## 2 Casimir Energy for 2-Tori.
In this section we compute the Casimir energy for various 6D fields in a compactification to 4 dimensions on a 2-torus $`𝒯_2`$, for a very broad class of boundary conditions for these fields about the two cycles of $`𝒯_2`$. Although the result is interesting in its own right, it is also the starting point for the later calculation of Casimir energy on orbifolds (which exhibit the effects for the Casimir energy of the presence of co-dimension 2 branes).
We define $`𝒯_2`$ by identifying points on the plane according to
$$(y_1,y_2)(y_1+n_2L_2\mathrm{cos}\theta +n_1L_1;y_2+n_2L_2\mathrm{sin}\theta ),$$
(1)
where $`n_{1,2}`$ are integers and $`\theta `$, $`L_1`$ and $`L_2`$ are the three real moduli of the torus (see Fig. 1). Equivalently, in terms of the complex coordinate $`z=y_1+iy_2`$ this is
$$zz+(n_2U+n_1)L_1,$$
(2)
with the complex quantity $`U`$ defined by $`U=\mathrm{exp}(i\theta )L_2/L_1=U_1+iU_2`$, $`(U_2>0)`$.
### 2.1 Scalar Field Casimir Energy
Consider a complex 6D field $`\mathrm{\Phi }`$ on a space-time compactified to 4 dimensions on the above 2-torus. Writing the six coordinates as $`\{x_\mu ,y_i\}`$, with $`\mu =0,\mathrm{},3`$ and $`i=1,2`$, assume the scalar satisfies the following boundary conditions
$`\mathrm{\Phi }(x,y_1+n_2L_2\mathrm{cos}\theta +n_1L_1;y_2+n_2L_2\mathrm{sin}\theta )=e^{2\pi i(n_1\rho _1+n_2\rho _2)}\mathrm{\Phi }(x,y_1,y_2),`$ (3)
with $`\rho _{1,2}`$ being two real quantities. The choices $`\rho _{1,2}=0,\frac{1}{2}`$ correspond to periodic or anti-periodic boundary conditions along the torus’ two cycles. More general values of $`\rho _i`$ are also possible, such as when $`\mathrm{\Phi }`$ transforms non-trivially under a gauge group for which nonzero Wilson lines are turned on for the corresponding cycles.<sup>1</sup><sup>1</sup>1To see this (for an orthogonal torus) notice that the toroidal boundary conditions preclude removing a constant gauge potential, such as $`A_1=a`$, using only strictly periodic gauge transformations. (This corresponds to the Wilson line $`W=_0^{L_1}A_m𝑑y^m=aL_1`$.) However, $`A_1`$ can be removed using a singular gauge transformation having parameter $`\omega =ay_1`$, at the expense of changing the boundary conditions of charged fields: $`\mathrm{\Phi }(y_1+L_1)=e^{iqaL_1}\mathrm{\Phi }(y_1)`$, where $`q`$ is $`\mathrm{\Phi }`$’s charge. We see from this that $`2\pi \rho _1=qaL_1=qW`$. Here we consider $`\rho _{1,2}`$ arbitrary.
Expanding the scalar field in terms of eigenfunctions of the 2D Laplacian, $`\mathrm{}_2=_1^2+_2^2`$, according to
$$\mathrm{\Phi }(x,y)=\underset{n_1,n_2}{}\varphi _{n_1,n_2}(x)f_{n_1,n_2}(y_1,y_2;\rho _1,\rho _2),$$
(4)
we have $`\mathrm{}_2f_{n_1,n_2}=M_{n_1,n_2}^2f_{n_1,n_2}`$ with
$`M_{n_1,n_2}^2(\rho _1,\rho _2)`$ $`=`$ $`{\displaystyle \frac{(2\pi )^2}{𝒜U_2}}|n_2+\rho _2U(n_1+\rho _1)|^2.`$ (5)
Here $`𝒜=L_1L_2\mathrm{sin}\theta `$ denotes the area of the torus. The mode functions are given by
$$f_{n_1,n_2}(y_1,y_2;\rho _1,\rho _2)=\frac{1}{\sqrt{𝒜}}e^{2i\pi [(n_1+\rho _1)(y_1y_2\mathrm{cot}\theta )/L_1+(n_2+\rho _2)y_2/(L_2\mathrm{sin}\theta )]}.$$
(6)
We now compute the vacuum energy for a complex scalar field having 6D mass $`m`$, compactified on $`𝒯_2`$ with the above boundary conditions. Denoting the vacuum-energy per unit 3-volume for such a scalar field by $`V(\rho _1,\rho _2)`$, we have
$`V(\rho _1,\rho _2)`$ $``$ $`\mu ^{4d}{\displaystyle \underset{n_{1,2}𝐙}{}}{\displaystyle \frac{d^dp}{(2\pi )^d}\mathrm{ln}\left[\frac{p^2+M_{n_1,n_2}^2+m^2}{\mu ^2}\right]}`$ (7)
$`=`$ $`{\displaystyle \frac{\mu ^4}{(2\pi )^d}}{\displaystyle \underset{n_{1,2}𝐙}{}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t^{1+d/2}}}e^{\pi t\left[(M_{n_1,n_2}^2+m^2)/\mu ^2\right]}`$
where we have continued the momentum integration to Euclidean signature, and we regulate the ultraviolet divergences which arise in the sum and integral using dimensional regularization, with the complex quantity<sup>2</sup><sup>2</sup>2The conventions here differ from those used in our previous ref. where $`d=42ϵ`$. $`d=4ϵ`$ ultimately being taken to 4. Here $`\mu `$ denotes the arbitrary mass scale which arises in dimensional regularization, and which drops out of all physical quantities.
A potential subtlety arises in the above expression for massless 6D fields ($`m=0`$) if both $`\rho _1`$ or $`\rho _2`$ are integers, because in this case there is a choice of integers $`(n_1,n_2)`$ for which $`M_{n_1,n_2}`$ vanishes. In this case it is convenient to keep $`m`$ nonzero so that all manipulations remain well-defined, with $`m`$ taken to zero at the end of the calculation. In principle, one must be alive to the possibility of unexpected singularities appearing when $`m`$ tends to zero after renormalization, such as the familiar infrared mass singularities of Quantum Electrodynamics . As usual, such infrared problems are less severe in higher dimensions and we shall see that there is no such obstruction to taking $`m0`$ for our applications in 6 dimensions.
The calculation of the two infinite sums in (7) is tedious, and is given in detail in Appendices A and Bc.f. eqs.(A-1), (B-1), (B) to (B-18) — using the approach of . In what follows we quote only the final results which are appropriate to the discussion at hand. The next three sections respectively concentrate on the ultraviolet-divergent part of the result, as well as the finite part in the cases where the 6D scalars are either massless ($`m^2𝒜0`$) or very massive ($`m^2𝒜1`$).
#### 2.1.1 Ultraviolet Divergences for 2-Tori
The ultraviolet divergent part of $`V`$ in (7) denoted $`V_{\mathrm{}}`$ is, with $`ϵ=4d`$ (see eqs. (B-1), (B) to (B-18))
$$V_{\mathrm{}}(\rho _1,\rho _2)=\frac{m^6𝒜}{192\pi ^3ϵ}$$
(8)
which is valid for arbitrary $`m`$. (Eq.(8) shows the importance of keeping a non-zero $`m`$ when discussing ultraviolet divergences in dimensional regularization (DR), since these can easily be missed if $`m=0`$.) This expression has several features on which we now remark (and which agree with the more general analysis of the ultraviolet divergences in 6D field theories compactified on Ricci-flat backgrounds given in a companion paper ).
$``$ First, the divergent part depends on the moduli, $`L_{1,2}`$ and $`\theta `$, only through the toroidal area $`𝒜=L_1L_2\mathrm{sin}\theta `$, and is interpreted as being a renormalization of the 6D cosmological constant. Note also that the UV divergence vanishes if we take $`\theta 0`$, corresponding to collapsing the two cycles of the torus onto each other down to one dimension less.<sup>3</sup><sup>3</sup>3This limit must be treated with care, however, since the calculations of Appendix A also require $`U_2\mathrm{sin}\theta `$ to be finite and nonzero. This agrees with the well-known absence of one-loop UV divergences for a broad class of theories when they are dimensionally regularized in odd dimensions.
$``$ The proportionality to $`m^6`$ is also what is required on dimensional grounds (in dimensional regularization) for a contribution to the 6D cosmological constant. The absence of other powers of $`m`$, such as an $`𝒜`$-independent result proportional to $`m^4`$, is a consequence of the torus being flat, and is not true for more general curved spacetimes . Once the UV divergence is renormalized into the 6D cosmological constant there is no obstruction to taking $`m0`$, unlike the situation for massless 4D theories. We consequently feel free to simply set $`m=0`$ in subsequent applications of our formulae for $`V^{\mathrm{ren}}`$ where
$`V^{\mathrm{ren}}VV_{\mathrm{}}`$ (9)
with $`V`$ as in (7) and $`V_{\mathrm{}}`$ its divergent part.
$``$ The divergent part of $`V`$ is $`\rho _i`$-independent and so does not depend on the boundary conditions of the scalar field, eqs.(3). Again this agrees with general arguments, since the short-wavelength modes responsible for the UV properties are not sensitive to the boundary conditions which depend on the global properties of the background geometry. We shall see that for orbifolds new divergences are present corresponding to counterterms localized at the fixed points, and these new divergences can depend on the nature of the boundary conditions imposed on the covering space.
$``$ The $`1/ϵ`$ pole which appears here represents a bona fide 6D divergence. This is at first sight surprising, since $`ϵ`$ represents the difference between $`d`$ and 4 rather than 6, and it is introduced for each of the Kaluza-Klein (KK) modes. To understand this it is important to recognize that our expressions contain two separate sources of UV divergence: the integration over 4-momentum, $`p`$, and the two sums over KK mode numbers, $`n_i`$. In our calculations the dimensional continuation is $`ϵ=4d`$ away from four in order to regularize the $`p`$-integration. On the other hand, the KK mode sums are managed using zeta-function techniques, and the presence of $`ϵ`$ ensures the regularisation of these sums. With these choices the leading inverse powers of $`ϵ`$ obtained turn out to be precisely those which would be obtained starting from 6D and following the powers of $`ϵ^{}=(6d)`$. This equivalence is shown in more detail for an explicit example (using a spherical geometry, for which more divergences may be followed) in ref. .
We now examine the finite parts of the Casimir energy density.
#### 2.1.2 Massless Fields in 6D
In the massless limit, $`m0`$, the result for the vacuum energy density is (see eqs.(A-3), (B-1), (B) to (B-18))
$`V^{\mathrm{ren}}(\rho _1,\rho _2)|_{m=0}={\displaystyle \frac{1}{𝒜^2}}\{{\displaystyle \frac{(2\pi U_2)^3}{90}}[{\displaystyle \frac{1}{21}}\mathrm{\Delta }_{\rho _1}^2(15\mathrm{\Delta }_{\rho _1}^22\mathrm{\Delta }_{\rho _1}^4+6\mathrm{\Delta }_{\rho _1}^3)]`$
$`+{\displaystyle \underset{n_1𝐙}{}}[(n_1+\mathrm{\Delta }_{\rho _1})^2\text{Li}_3(\sigma _{n_1})+{\displaystyle \frac{3\left|n_1+\mathrm{\Delta }_{\rho _1}\right|}{2\pi U_2}}\text{Li}_4(\sigma _{n_1})+{\displaystyle \frac{3}{(2\pi U_2)^2}}\text{Li}_5(\sigma _{n_1})+c.c.]\}`$
with
$`\sigma _{n_10}=e^{2i\pi \left(\mathrm{\Delta }_{\rho _2}U\mathrm{\Delta }_{\rho _1}Un_1\right)}\text{and}\sigma _{n_1<0}={\displaystyle \frac{1}{\overline{\sigma }_{n_1>0}}},`$ (10)
where $`\overline{\sigma }_{n_1}`$ denotes the complex conjugate of $`\sigma _{n_1}`$. In these expressions $`0\mathrm{\Delta }_{\rho _i}<1`$ represents the fractional part of $`\rho _i`$, as in $`\rho _i=[\rho _i]+\mathrm{\Delta }_{\rho _i}`$, where $`[\rho _i]𝐙`$ is the largest integer smaller than or equal to $`\rho _i`$. The poly-logarithm functions which appear here are defined by the sums
$$\text{Li}_\sigma (x)=\underset{n=1}{\overset{\mathrm{}}{}}\frac{x^n}{n^\sigma }.$$
(11)
The point of rewriting the initial two sums into the ones written here is that these converge well and so are useful for numerical purposes. Figure 2 plots $`V(\rho _1,\rho _2)`$ as functions of the moduli $`U_1`$ and $`U_2`$ for various choices for the boundary conditions $`(\rho _1,\rho _2)`$.
We have checked that the above formula agrees with the particular cases studied in the literature. For instance in the special case $`\rho _1=\rho _2=0`$ we find
$`V^{\mathrm{ren}}(0,0)|_{m=0}`$ $`=`$ $`{\displaystyle \frac{1}{𝒜^2}}\{{\displaystyle \frac{4\pi ^3U_2^3}{945}}+{\displaystyle \frac{3\zeta [5]}{2\pi ^2U_2^2}}+2{\displaystyle \underset{n_1=1}{\overset{\mathrm{}}{}}}[n_1^2\text{Li}_3(q^{n_1})`$ (12)
$`+{\displaystyle \frac{3n_1}{2\pi U_2}}\text{Li}_4(q^{n_1})+{\displaystyle \frac{3}{4\pi ^2U_2^2}}\text{Li}_5(q^{n_1})+c.c.]\},`$
where $`qe^{2i\pi U}`$. This agrees with the result given in ref. .
#### 2.1.3 Heavy-Mass Dependence
The generality of the calculation in Appendix A, B also allows the explicit exhibition of the heavy-mass limit, $`m^2𝒜\mathrm{}`$, of the Casimir energy. This is of particular interest for Supersymmetric Large Extra Dimensions, where the naturalness of the description of the Dark Energy density relies on the Casimir energy only depending weakly on the masses of heavy fields in the 6D bulk. Using formulae (B), (B-16), (B-18) of the appendix it may be shown that if $`m\{1/L_1,1/(L_2\mathrm{sin}\theta )\}`$, leading to $`m^2𝒜1`$, then
$$V(\rho _1,\rho _2)=\frac{m^6𝒜}{384\pi ^3}\left(\frac{2}{ϵ}\mathrm{ln}\frac{4\pi ^3m^2e^{\gamma 11/6}}{\mu ^2}\right)+\frac{1}{𝒜^2}𝒪\left((m^2𝒜)^\mu U_2^\nu e^{2\pi (m^2𝒜)^\sigma U_2^{\pm \sigma }}\right)$$
(13)
thus powers of $`m^2𝒜`$ other than $`m^6`$ are exponentially suppressed. These expressions agree well with the general results of ref. , which identify the large-$`m`$ behaviour using general heat-kernel techniques. For general geometries there can be powers of $`m`$ in the large-$`m`$ limit, but the leading such powers are proportional to local effective interactions which involve polynomials of the background fields and their derivatives. For the simple toroidal geometries considered here all of these local interactions vanish, leading to the exponential mass suppression found above.
The absence of powers like $`m^4`$ or $`m^2`$ in the large-$`m`$ limit is more difficult to understand from the point of view where the 6D calculation is regarded as simply being the sum over an infinite number of 4D contributions, each of which can themselves have such powers of $`m`$. As is clear from the general 6D analysis of , the absence of these terms may be traced to the requirements of locality and general covariance in 6 dimensions — requirements which are easily missed in a KK mode sum calculation.
### 2.2 Higher-Spin Fields on $`𝒯_2`$.
With an eye towards applications to supersymmetric theories, in this section we compute the corresponding results for the Casimir energy for other massless fields in 6 dimensions. We do so using the trick of ref. , which uses the prior knowledge that the Casimir energy must vanish once summed over the field content of a 6D supermultiplet, provided that these fields all share the same boundary conditions about the cycles of $`𝒯_2`$ (and so do not break any of the supersymmetries).
To this end we reproduce as Table 1 a table from ref. listing the massless field content of some of the representations of $`(2,0)`$ supersymmetry in 6 dimensions. In this table the scalars are real, the spinors are symplectic-Weyl and the 2-form gauge potentials are self-dual or anti-self-dual.
The argument of ref. uses the observation that a single symplectic-Weyl fermion and two real scalars preserve 6-dimensional $`(2,0)`$ supersymmetry in a toroidal compactification for which they share the same boundary conditions, $`(\rho _1,\rho _2)`$, about the torus’ two cycles. The Casimir energy for these fields must therefore cancel in order to give a vanishing result for the contribution of a hypermultiplet. Since we know the scalar result for general $`\rho _i`$, we may infer from this that the Casimir energy for a single symplectic-Weyl fermion must be precisely $`1`$ times the result quoted above for a complex scalar field having the same boundary conditions.
Using the identical argument based on the vanishing of the Casimir energy summed over the field content of a gauge multiplet similarly shows that the Casimir energy of a 6D gauge boson must be $`2`$ times the result for a 6D symplectic-Weyl fermion — and so is $`+2`$ times the result for a 6D complex scalar — having the same boundary conditions. Arguing in this way allows the inference of the Casimir energy for all of the other fields appearing in the supermultiplets listed in Table 1. The results found in this way are
$`V_{1/2}(\rho _1,\rho _2)`$ $`=`$ $`V(\rho _1,\rho _2)`$
$`V_1(\rho _1,\rho _2)`$ $`=`$ $`2V(\rho _1,\rho _2)`$
$`V_{KR}(\rho _1,\rho _2)`$ $`=`$ $`{\displaystyle \frac{3}{2}}V(\rho _1,\rho _2)`$
$`V_{3/2}(\rho _1,\rho _2)`$ $`=`$ $`3V(\rho _1,\rho _2)`$
$`V_2(\rho _1,\rho _2)`$ $`=`$ $`V(\rho _1,\rho _2),`$ (14)
where the divergent and finite parts of the right-hand side of this equation are given explicitly by eqs. (8) and (2.1.2), above. Here $`V_{1/2}`$, $`V_1`$, $`V_{KR}`$, $`V_{3/2}`$ and $`V_2`$ respectively denote the results for a symplectic-Weyl fermion, a gauge boson, a Kalb-Ramond self-dual (or anti-self-dual) 2-form gauge potential, a symplectic-Weyl gravitino and a graviton.
Given these expressions, it is simple to compute the nonzero Casimir energy which results when 6D supersymmetry is broken à la Scherk and Schwarz , by assigning different boundary conditions to different fields within a single supermultiplet. Supersymmetry is broken in this case by the boundary conditions themselves. For instance, if the symplectic-Weyl fermion in a hypermultiplet has boundary condition $`(\rho _1,\rho _2)`$ but the complex scalar has boundary condition $`(\rho _1^{},\rho _2^{})`$ then the Casimir energy for this hypermultiplet would be
$$V_{\mathrm{hyper}}=V_{1/2}(\rho _1,\rho _2)+V(\rho _1^{},\rho _2^{})=V(\rho _1^{},\rho _2^{})V(\rho _1,\rho _2).$$
(15)
Similarly applying different boundary conditions to the constituents of a gauge or tensor multiplet gives
$`V_{\mathrm{gauge}}`$ $`=`$ $`V_1(\rho _1,\rho _2)+2V_{1/2}(\rho _1^{},\rho _2^{})=2[V(\rho _1,\rho _2)V(\rho _1^{},\rho _2^{})]`$
$`V_{\mathrm{tensor}}`$ $`=`$ $`V_{KR}(\rho _1,\rho _2)+2V_{1/2}(\rho _1^{},\rho _2^{})+{\displaystyle \frac{1}{2}}V(\rho _1^{\prime \prime },\rho _2^{\prime \prime })`$ (16)
$`=`$ $`{\displaystyle \frac{3}{2}}V(\rho _1,\rho _2)2V(\rho _1^{},\rho _2^{})+{\displaystyle \frac{1}{2}}V(\rho _1^{\prime \prime },\rho _2^{\prime \prime }),`$
and so on.
Figure 2 gives in addition to the plots of $`V(\rho _1,\rho _2)`$, the differences $`V(\rho _1,\rho _2)V(\rho _1^{},\rho _2^{})`$ as functions of $`U_1`$ and $`U_2`$ for various choices for the boundary conditions $`(\rho _1,\rho _2)`$ and $`(\rho _1^{},\rho _2^{})`$. The periodicity wrt $`U_1`$ in the plots is a remnant of the $`SL(2,Z)_U`$ symmetry (modified by non-zero Wilson lines). For $`U_2=L_2/L_1\mathrm{sin}\theta 𝒪(1)`$, one has flat directions for $`V`$ (as function of $`U_1,U_2`$). This changes for $`U_21`$ (say if $`\theta 1`$) when $`V`$ develops maxima/minima. For values of $`\rho _i`$ other than those in the figure, the peaks in these plots have different height.
## 3 Casimir Energy for Orbifolds
None of the previous results included the effects of 3-branes within the bulk when computing the Casimir energies. We now extend these calculations to some simple examples which include branes, and for which the background geometry includes the back-reaction of the brane tensions by incorporating the appropriate conical singularities at the brane positions. We only consider here brane singularities which correspond to the specific defect angles which arise when an orbifold is constructed from the 2-torus by identifying points under the action, $`Z_N`$, of a discrete set of rotations. We analyze separately the case of $`𝒯_2/Z_2`$, with full details, and then the orbifold $`𝒯_2/Z_4`$ whose technical details differ considerably. For $`𝒯_2/Z_4`$ we use a general method which can be applied to the remaining $`𝒯_2/Z_3`$ and $`𝒯_2/Z_6`$.
To this end we return to the description of $`𝒯_2`$ as a complex plane, $`z=y_1+iy_2`$, identified under the action of a lattice of discrete translations as in eq. (2). Following standard practice, we construct an orbifold from this torus by further identifying points under the action of the $`Z_N`$ rotations defined by
$$z\tau ^kz,$$
(17)
where $`\tau =e^{2\pi i/N}`$. This gives a well-defined coset space, $`𝒪=𝒯_2/Z_N`$, provided that these rotations take the initial lattice which defines the torus onto itself. Notice that if $`N=2`$, then the rotation $`zz`$ is automatically a symmetry of the lattice for any value of the moduli $`L_1`$, $`L_2`$ and $`\theta `$, and so these three quantities are also moduli of the resulting orbifold, $`𝒪`$. On the other hand, if $`N>2`$ then the rotation is a symmetry of the lattice only for specific choices for the complex structure: $`U(L_2/L_1)e^{i\theta }=\tau `$, and so $`L_2=L_1=L`$ and $`\theta =2\pi /N`$, and so only one modulus, $`L`$, in this case survives.
The coset $`𝒯_2/Z_N`$ is an orbifold rather than a manifold because of the metric singularities which arise at the fixed points of the group. For instance, in the case $`𝒯_2/Z_2`$ there are 4 such points, corresponding to $`z=0,\frac{1}{2},\frac{1}{2}U`$ and $`\frac{1}{2}(1+U)`$. The metric has a conical singularity at each of these points, whose defect angle is $`\pi `$. For further details on orbifolds and their fixed points see Appendix D.
### 3.1 A Scalar Field on $`𝒯_2/Z_2`$.
We now compute the Casimir energy for a complex scalar field, $`\mathrm{\Phi }`$, compactified from 6D to 4D on the orbifold $`𝒪=𝒯_2/Z_2`$. As before, the 6 coordinates are taken to be $`\{x,y_i\}`$, with the orbifold corresponding to the coordinates $`y_i,i=1,2`$. We consider the 6D scalar field to satisfy the boundary conditions
$`\mathrm{\Phi }(x,y_1+L_1;y_2)`$ $`=`$ $`e^{2\pi i\rho _1}\mathrm{\Phi }(x,y_1,y_2);\rho _1=0,1/2.`$ (18)
$`\mathrm{\Phi }(x,y_1+L_2\mathrm{cos}\theta ;y_2+L_2\mathrm{sin}\theta )`$ $`=`$ $`e^{2\pi i\rho _2}\mathrm{\Phi }(x,y_1,y_2);\rho _2=0,1/2.`$ (19)
$`\mathrm{\Phi }(x,y_1,y_2)`$ $`=`$ $`\pm \mathrm{\Phi }(x,y_1,y_2).`$ (20)
Condition (20) is possible because of the new cycle that the orbifold has (which the torus does not). As is indicated in eqs. (18) and (19), the quantities $`\rho _i`$ are no longer free to take any real value in this case, because the underlying Wilson lines must be compatible with the orbifold rotation. For instance, when acting on the coordinates it is straightforward to show that the composite transformation $`X=\mathrm{\Sigma }P\mathrm{\Sigma }`$ gives $`X(y_1,y_2)=(y_1,y_2)`$, if $`\mathrm{\Sigma }(y_1,y_2)=(y_1+L_1,y_2)`$ and $`P(y_1,y_2)=(y_1,y_2)`$. Applying the same transformations to $`\mathrm{\Phi }`$ and using the boundary conditions (18) and (20), consistency requires $`\mathrm{exp}(4\pi i\rho _1)=1`$, and so $`\rho _1=n/2`$ for some integer $`n`$. A similar argument implies that $`\rho _2`$ is also half-integer. If we require, without loss of generality, $`0\rho _i<1`$, then we see that consistency requires $`\rho _{1,2}=0,1/2`$. Altogether there are 8 possible choices for the boundary conditions, denoted by $`(\rho _1,\rho _2)^\pm `$.
Because the orbifold reflection is a symmetry of $`\mathrm{}_2`$, these reflections have a natural action on its eigenfunctions, $`f_{n_1,n_2}`$. Using the explicit expressions obtained earlier, eq. (6), for the toroidal mode functions
$`f_{n_1,n_2}(y_1,y_2;\rho _1,\rho _2)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{𝒜}}}e^{2i\pi [(n_1+\rho _1)(y_1y_2\mathrm{cot}\theta )/L_1+(n_2+\rho _2)y_2/(L_2\mathrm{sin}\theta )]}`$ (21)
$`=`$ $`f_{n_1,n_2}^{}(y_1,y_2;\rho _1,\rho _2),`$
$``$ $`f_{n_1^{},n_2^{}}(y_1,y_2;\rho _1,\rho _2),`$
where $`n_i^{}=n_i2\rho _i`$ ensures $`n_i^{}+\rho _i=(n_i+\rho _i)`$. Notice that $`n_i^{}`$ defined in this way remains an integer because for the $`Z_2`$ orbifold $`\rho _i=0,\frac{1}{2}`$.
Using this action it is straightforward to specialize the toroidal mode expansion, eq. (4), to fields on $`𝒯_2/Z_2`$ with boundary conditions $`(\rho _1,\rho _2)^\pm `$. We now write these expansions explicitly, in terms of the real and imaginary parts of the mode functions $`f_{n_1,n_2}=e_{n_1,n_2}+ig_{n_1,n_2}`$.
(1). $`(\rho _1,\rho _2)=(0,0)^\pm `$. In this case $`f_{n_1,n_2}(y_1,y_2)=f_{n_1,n_2}(y_1,y_2)`$ and so using the boundary conditions of eqs. (20) in the mode expansion of eq. (4) gives
$`\mathrm{\Phi }^+(x,y_1,y_2)`$ $`=`$ $`{\displaystyle \underset{n_10,n_2𝐙}{}}{\displaystyle \frac{2}{2^{\delta _{n_1,0}}}}\varphi _{n_1,n_2}(x)e_{n_1,n_2}(y_1,y_2;0,0),`$
$`\mathrm{\Phi }^{}(x,y_1,y_2)`$ $`=`$ $`{\displaystyle \underset{n_10,n_2𝐙}{\overset{}{}}}{\displaystyle \frac{2i}{2^{\delta _{n_1,0}}}}\varphi _{n_1,n_2}(x)g_{n_1,n_2}(y_1,y_2;0,0),`$ (22)
where the superscript on $`\mathrm{\Phi }`$ indicates the sign chosen for the orbifold projection, and $`\delta _{n_1,0}`$ is the usual Kronecker delta-function which vanishes unless $`n_1=0`$, in which case it equals unity. Notice that because $`g_{0,0}(y_1,y_2;0,0)=0`$, the mode $`(n_1,n_2)=(0,0)`$ is absent in the double sum for $`\mathrm{\Phi }^{}`$, as is indicated by the primed double sum.
(2). $`(\rho _1,\rho _2)=(0,1/2)^\pm `$. In this case $`f_{n_1,n_2}(y_1,y_2)=f_{n_1,n_21}(y_1,y_2)`$ and so the mode expansion becomes
$`\mathrm{\Phi }^+(x,y_1,y_2)`$ $`=`$ $`{\displaystyle \underset{n_10,n_2𝐙}{}}{\displaystyle \frac{2}{2^{\delta _{n_1,0}}}}\varphi _{n_1,n_2}(x)e_{n_1,n_2}(y_1,y_2;0,1/2),`$
$`\mathrm{\Phi }^{}(x,y_1,y_2)`$ $`=`$ $`{\displaystyle \underset{n_10,n_2𝐙}{}}{\displaystyle \frac{2i}{2^{\delta _{n_1,0}}}}\varphi _{n_1,n_2}(x)g_{n_1,n_2}(y_1,y_2;0,1/2).`$ (23)
(3). $`(\rho _1,\rho _2)=(1/2,0)^\pm `$. Here $`f_{n_1,n_2}(y_1,y_2)=f_{n_11,n_2}(y_1,y_2)`$, and so
$`\mathrm{\Phi }^+(x,y_1,y_2)`$ $`=`$ $`{\displaystyle \underset{n_10,n_2𝐙}{}}2\varphi _{n_1,n_2}(x)e_{n_1,n_2}(y_1,y_2;1/2,0),`$
$`\mathrm{\Phi }^{}(x,y_1,y_2)`$ $`=`$ $`{\displaystyle \underset{n_10,n_2𝐙}{}}2i\varphi _{n_1,n_2}(x)g_{n_1,n_2}(y_1,y_2;1/2,0),`$ (24)
(4). $`(\rho _1,\rho _2)=(1/2,1/2)^\pm `$. In this case $`f_{n_1,n_2}(y_1,y_2)=f_{n_11,n_21}(y_1,y_2)`$, and so
$`\mathrm{\Phi }^+(x,y_1,y_2)`$ $`=`$ $`{\displaystyle \underset{n_10,n_2𝐙}{}}2\varphi _{n_1,n_2}(x)e_{n_1,n_2}(y_1,y_2;1/2,1/2),`$
$`\mathrm{\Phi }^{}(x,y_1,y_2)`$ $`=`$ $`{\displaystyle \underset{n_10,n_2𝐙}{}}2i\varphi _{n_1,n_2}(x)g_{n_1,n_2}(y_1,y_2;1/2,1/2).`$ (25)
There are two ways to compute the Casimir energy for the scalar field $`\mathrm{\Phi }`$ on $`𝒯_2/Z_2`$ with these boundary conditions. One approach is to recognize that the scalar propagator on the orbifold may be obtained from the propagator on the torus using the method of images:
$$G_𝒪^\pm (xx^{},yy^{})=G_𝒯(xx^{},yy^{})\pm G_𝒯(xx^{},y+y^{}),$$
(26)
and following the implications of this for the vacuum energy. The second approach is to directly perform the KK mode sum over the modified mode functions given above. Both lead to the same result, and we present the mode-function derivation here because, albeit more involved, it can be extended to the case of $`𝒯_2/Z_N`$ and allows a general discussion of the ultraviolet divergences for $`𝒯_2/Z_N`$.
The vacuum energy density per unit 3-volume written as a mode sum is given by
$`V_𝒪(0,0)^\pm `$ $`=`$ $`{\displaystyle \underset{n_10,n_2𝐙}{\overset{()^{}}{}}}{\displaystyle \frac{\mu ^{4d}}{2^{\delta _{n_1,0}}}}{\displaystyle \frac{d^dp}{(2\pi )^d}\mathrm{ln}\left[\frac{p^2+M_{n_1,n_2}^2(0,0)+m^2}{\mu ^2}\right]}`$
$`V_𝒪(0,1/2)^\pm `$ $`=`$ $`{\displaystyle \underset{n_10,n_2𝐙}{}}{\displaystyle \frac{\mu ^{4d}}{2^{\delta _{n_1,0}}}}{\displaystyle \frac{d^dp}{(2\pi )^d}\mathrm{ln}\left[\frac{p^2+M_{n_1,n_2}^2(0,1/2)+m^2}{\mu ^2}\right]}`$
$`V_𝒪(1/2,0)^\pm `$ $`=`$ $`{\displaystyle \underset{n_10,n_2𝐙}{}}\mu ^{4d}{\displaystyle \frac{d^dp}{(2\pi )^d}\mathrm{ln}\left[\frac{p^2+M_{n_1,n_2}^2(1/2,0)+m^2}{\mu ^2}\right]}`$
$`V_𝒪(1/2,1/2)^\pm `$ $`=`$ $`{\displaystyle \underset{n_10,n_2𝐙}{}}\mu ^{4d}{\displaystyle \frac{d^dp}{(2\pi )^d}\mathrm{ln}\left[\frac{p^2+M_{n_1,n_2}^2(1/2,1/2)+m^2}{\mu ^2}\right]},`$ (27)
where $`d=(4ϵ)`$ and the symbol $`()^{}`$ on the sum in the first line indicates the exclusion from the sum of the single mode $`(n_1,n_2)=(0,0)`$, but only for the case of $`(0,0)^{}`$ boundary conditions.
As might be expected from the approach based on the method of images, these expressions may be evaluated in terms of the corresponding quantities on the torus. To see this for the mode sums we denote the summands of these expressions by $`𝒲(|n_2+\rho _2U(n_1+\rho _1)|^2)`$ in order to emphasize their dependence on mode numbers and moduli. It is then simple to use the invariance of $`𝒲`$ under changes in sign of $`(n_i+\rho _i)`$ to prove the following identities:
$`{\displaystyle \underset{n_10;n_2𝐙}{}}{\displaystyle \frac{1}{2^{\delta _{n_1,0}}}}𝒲\left(|n_2Un_1|^2\right)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{n_1,n_2𝐙}{}}𝒲\left(|n_2Un_1|^2\right)`$
$`{\displaystyle \underset{n_10;n_2𝐙}{}}{\displaystyle \frac{1}{2^{\delta _{n_1,0}}}}𝒲\left(|n_2+1/2Un_1|^2\right)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{n_1,n_2𝐙}{}}𝒲\left(|n_2+1/2Un_1|^2\right)`$
$`{\displaystyle \underset{n_10;n_2𝐙}{}}𝒲\left(|n_2+\rho _2U(n_1+1/2)|^2\right)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{n_1,n_2𝐙}{}}𝒲\left(|n_2+\rho _2U(n_1+1/2)|^2\right),`$ (28)
where $`\rho _2`$ is either 0 or $`\frac{1}{2}`$ in the last line. These expressions allow the derivation of the following expressions for the Casimir energies in terms of the toroidal results, $`V(\rho _1,\rho _2)`$
$`V_𝒪(0,0)^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[V(0,0)V_{\mathrm{zm}}\right]`$
$`V_𝒪(\rho _1,\rho _2)^\pm `$ $`=`$ $`{\displaystyle \frac{1}{2}}V(\rho _1,\rho _2)\text{for all others.}`$ (29)
Here $`V_{\mathrm{zm}}`$ is the contribution to the torus Casimir energy of the “zero mode” $`(n_1,n_2)=(0,0)`$ with $`(\rho _1,\rho _2)^{}=(0,0)^{}`$ (for its expression see (B-17) and (B-18)). We can now present in detail the divergent and finite parts of the sums and integrals in (3.1), (3.1).
#### 3.1.1 Ultraviolet Divergences
We isolate the divergent part of $`V_𝒪`$ in eqs.(3.1) and write
$`V_𝒪(\rho _1,\rho _2)^\pm =V_{𝒪,\mathrm{}}(\rho _1,\rho _2)^\pm +V_𝒪^{\mathrm{ren}}(\rho _1,\rho _2)^\pm `$ (30)
where all divergent terms are included in $`V_{𝒪,\mathrm{}}`$. As is clear from eqs. (3.1), for all choices of boundary condition except $`(0,0)^{}`$ on the orbifold the ultraviolet divergences encountered are precisely half of those encountered on the torus, eq. (8):
$`V_𝒪\mathrm{}(\rho _1,\rho _2)^\pm ={\displaystyle \frac{m^6𝒜}{384\pi ^3ϵ}}={\displaystyle \frac{m^6𝒜_𝒪}{192\pi ^3ϵ}},\text{if}(\rho _1,\rho _2)^\pm (0,0)^{}.`$ (31)
This divergence may be absorbed, as usual, into a renormalization of the bulk cosmological constant. Notice that its coefficient is the same as was obtained earlier for the torus, once the divergence is expressed in terms of the area of the orbifold, $`𝒜_𝒪`$, which is half the area, $`𝒜`$, of the covering torus.
By contrast, the exclusion of the zero mode for the specific choice $`(0,0)^{}`$ introduces a new type of divergence which was not encountered for the torus. In this case the orbifold and toroidal divergences differ by the contribution of the $`n_1=n_2=0`$ mode alone, and thus, using eqs.(A-3), (B-1), (B), (B), one has
$`V_𝒪\mathrm{}(0,0)^{}={\displaystyle \frac{1}{2}}V_{\mathrm{}}^{}(0,0)={\displaystyle \frac{m^6𝒜_𝒪}{192\pi ^3ϵ}}+{\displaystyle \frac{m^4}{32\pi ^2ϵ}}.`$ (32)
The presence of the last term is consistent with the general heat-kernel analysis. Because it is proportional to $`m^4`$ and is independent of the bulk moduli, it has the right properties to be interpreted as a renormalization of the tension of the branes whose presence at the fixed points is responsible for the conical singularities in the bulk geometry at these points. As before, the ultraviolet-finite part of the Casimir energy obtained after this renormalization is nonsingular in the $`m0`$ limit, and so we are free to take this limit explicitly in the renormalized result for massless 6D fields which we quote below.
It should be emphasized that although this divergence renormalizes the local brane tensions, it arises due to the functional integration over bulk fields. This is a feature which arises quite generically for quantum effects in the presence of boundaries and defects, whose origin can be understood in detail as follows. The bulk vacuum energy density, $`t_{MN}=T_{MN}`$ is ultraviolet finite (for the flat orbifold under discussion) after the bulk cosmological constant is appropriately renormalized. However although $`t_{MN}`$ is finite, it is also position-dependent due to the presence of the orbifold singularities breaking the translation invariance of the underlying torus. In particular, $`t_{MN}(y)`$ typically goes to infinity as the singular points are approached in a way which diverges once integrated over the volume of the orbifold. It is this new divergence which is renormalized by the brane-tension counter-term localized at the singularity.
#### 3.1.2 6D Massless Fields
Using eqs.(3.1) we can now give the explicit results for the Casimir energy of a massless complex 6D scalar field compactified on $`𝒯_2/Z_2`$, for the various boundary conditions $`(\rho _1,\rho _2)^\pm `$. It is noteworthy that $`V_{\mathrm{zm}}^{\mathrm{ren}}(0,0)`$ vanishes as $`m0`$ in dimensional regularization, and so the orbifold result is half of the appropriate toroidal result for all choices of boundary conditions. After the renormalization of the ultraviolet divergences described above, one has
$`V_𝒪^{\mathrm{ren}}(\rho _1,\rho _2)^\pm |_{m=0}={\displaystyle \frac{1}{2}}V^{\mathrm{ren}}(\rho _1,\rho _2)|_{m=0},\rho _{1,2}=0,1/2.`$ (33)
and that
$`V_𝒪^{\mathrm{ren}}(0,0)^\pm |_{m=0}={\displaystyle \frac{1}{𝒜^2}}\{\left({\displaystyle \frac{1}{21}}\right){\displaystyle \frac{(2\pi U_2)^3}{180}}+{\displaystyle \frac{3\zeta [5]}{(2\pi U_2)^2}}`$
$`+{\displaystyle \underset{n_1=1}{\overset{\mathrm{}}{}}}[n_1^2\text{Li}_3(q^{n_1})+{\displaystyle \frac{3n_1}{2\pi U_2}}\text{Li}_4(q^{n_1})+{\displaystyle \frac{3}{(2\pi U_2)^2}}\text{Li}_5(q^{n_1})+c.c.]\}`$
$`V_𝒪^{\mathrm{ren}}(0,1/2)^\pm |_{m=0}={\displaystyle \frac{1}{𝒜^2}}\{\left({\displaystyle \frac{1}{21}}\right){\displaystyle \frac{(2\pi U_2)^3}{180}}{\displaystyle \frac{45\zeta [5]}{16(2\pi U_2)^2}}`$
$`+{\displaystyle \underset{n_1=1}{\overset{\mathrm{}}{}}}[n_1^2\text{Li}_3(q^{n_1})+{\displaystyle \frac{3n_1}{2\pi U_2}}\text{Li}_4(q^{n_1})+{\displaystyle \frac{3}{(2\pi U_2)^2}}\text{Li}_5(q^{n_1})+c.c.]\}`$
$`V_𝒪^{\mathrm{ren}}(1/2,0)^\pm |_{m=0}={\displaystyle \frac{1}{𝒜^2}}\{\left({\displaystyle \frac{31}{672}}\right){\displaystyle \frac{(2\pi U_2)^3}{180}}+{\displaystyle \underset{n_1=0}{\overset{\mathrm{}}{}}}[(n_1+1/2)^2\text{Li}_3\left(q^{n_1+1/2}\right)`$
$`+{\displaystyle \frac{3\left(n_1+1/2\right)}{2\pi U_2}}\text{Li}_4\left(q^{n_1+1/2}\right)+{\displaystyle \frac{3}{(2\pi U_2)^2}}\text{Li}_5\left(q^{n_1+1/2}\right)+c.c.]\}`$
$`V_𝒪^{\mathrm{ren}}(1/2,1/2)|_{m=0}={\displaystyle \frac{1}{𝒜^2}}\{\left({\displaystyle \frac{31}{672}}\right){\displaystyle \frac{(2\pi U_2)^3}{180}}+{\displaystyle \underset{n_1=0}{\overset{\mathrm{}}{}}}[(n_1+1/2)^2\text{Li}_3(q^{n_1+1/2})`$
$`+{\displaystyle \frac{3\left(n_1+1/2\right)}{2\pi U_2}}\text{Li}_4(q^{n_1+1/2})+{\displaystyle \frac{3}{(2\pi U_2)^2}}\text{Li}_5(q^{n_1+1/2})+c.c.]\}`$ (34)
where $`q=e^{2\pi iU}`$ and the complex conjugate applies only to the series of polylogarithms.
#### 3.1.3 Heavy-Mass Dependence
The divergences of the Casimir energy for large $`m`$ are identical to those for the case of small $`m`$ discussed in Section 3.1.1. Further, because the orbifold results are simply expressed in terms of the toroidal ones, the heavy-mass dependence of the toroidal expressions carry over immediately to the orbifold Casimir energy. In particular, in dimensional regularization (and after modified minimal subtraction) the finite parts of the Casimir energy fall exponentially for large $`m`$, and the only strong $`m`$-dependence arises in the divergent terms, including the new $`m^4`$ term which arises for some of the boundary conditions.
#### 3.1.4 Higher-Spin Fields on $`𝒯_2/Z_2`$
The results for massless higher-spin fields on the $`𝒯_2/Z_2`$ orbifold can be read from their toroidal counterparts of Section 2.2. This is possible because the orbifold identification does not break supersymmetry provided that all of the fields within a 6D supermultiplet satisfy the same boundary conditions. The results for the Casimir energy of higher-spin fields may therefore simply be read off by multiplying the expressions (3.1.2) by the factors given in eqs. (2.2).
### 3.2 The Orbifold $`𝒯_2/Z_N`$ with $`N>2`$.
In this section we outline the steps for computing the Casimir energy for a complex scalar field $`\mathrm{\Phi }`$ compactified on the $`𝒯_2/Z_N`$ orbifolds, with $`N>2`$. Recall that for these orbifolds
$`U=e^{2i\pi /N},L_2=L_1=L,\theta =2\pi /N`$ (35)
We take the following action of the translation and orbifold $`Z_N`$ symmetries on the field $`\mathrm{\Phi }`$
$`\mathrm{\Phi }^g(x,\tau ^kz)`$ $`=`$ $`g^k\mathrm{\Phi }^g(x,z),`$
$`\mathrm{\Phi }^g(x,z+L_1)`$ $`=`$ $`e^{2i\pi \rho _1}\mathrm{\Phi }^g(x,z)`$
$`\mathrm{\Phi }^g(x,z+UL_1)`$ $`=`$ $`e^{2i\pi \rho _2}\mathrm{\Phi }^g(x,z),`$ (36)
where we use complex coordinates $`z=y_1+iy_2`$ and as before $`\tau e^{2i\pi /N}`$ and $`k=\overline{0,N1}`$. Here $`g`$ is a particular representation of the $`Z_N`$ transformation, acting on $`\mathrm{\Phi }^g`$. The superscript ‘$`g`$’ emphasizes that this representation is not unique, and the explicit form taken by $`\mathrm{\Phi }^g`$ in general depends on which $`g`$ is chosen. As in the case of $`𝒯_2/Z_2`$, this realization only faithfully reproduces the symmetry for specific choices for the $`\rho _i`$, whose values we now determine.
The consistency conditions for the action on $`\mathrm{\Phi }^g`$ are found by combining the above expressions and using geometrical relations which state how some of the $`Z_N`$ rotations can also be expressed as translations on the covering torus. To display these we use the complex coordinate $`z=y_1+iy_2`$, in terms of which the lattice of translations which defines the underlying torus is generated by $`e_1=L_1`$ and $`e_2=UL_1`$. Then, depending on the group of rotations, $`Z_N`$, which is of interest, the following restrictions can arise.
* If an orbifold rotation takes $`e_1`$ to $`e_2`$i.e. there is an integer $`0<k<N`$ for which $`\tau ^ke_1=e_2`$ (or $`\tau ^k=U`$) — then using eqs. (3.2) to evaluate $`\mathrm{\Phi }^g(x,\tau ^k(z+e_1))=\mathrm{\Phi }^g(x,\tau ^kz+e_2)`$ implies
$$g^k\mathrm{exp}(2\pi i\rho _1)=g^k\mathrm{exp}(2\pi i\rho _2),$$
(37)
and so we may take $`\rho _1=\rho _2`$ without loss of generality.
* If an orbifold rotation takes $`e_ie_i`$ then a similar argument implies
$$g^k\mathrm{exp}(2\pi i\rho _i)=g^k\mathrm{exp}(2\pi i\rho _i),$$
(38)
and so we may take $`\rho _i=0`$ or $`1/2`$.
For instance, only the second of these conditions applied to the $`Z_2`$ orbifold considered previously. By contrast, both conditions apply to the case of $`\pi /N`$ rotations which give the orbifold $`𝒯_2/Z_{2N}`$, and so for this case we must take $`\rho _1=\rho _2=0,\frac{1}{2}`$. For the $`Z_3`$ case, on the other hand, the first condition applies but instead of the second condition one has $`e_2+\tau ^ke_2=e_1`$, and so we find $`\rho _2=\rho _1=0,\frac{1}{3}`$ or $`\frac{2}{3}`$. These results express the quantization on these orbifolds of the underlying Wilson lines which are responsible for the boundary conditions which are expressed by the $`\rho _i`$. For a more detailed description of Wilson lines and their values on orbifolds see Appendix D.
To determine the action of the symmetries (3.2) on the toroidal mode functions, we adapt the discussion of Section 2.3 of ref. to include the general phases $`\rho _{1,2}`$. It is convenient for these purposes to rewrite eq. (6) in complex coordinates
$$f_{\underset{¯}{n}_1,\underset{¯}{n}_2}(z,\overline{z})=\frac{1}{\sqrt{𝒜}}e^{[(\underset{¯}{n}_2\underset{¯}{n}_1\overline{U})z(\underset{¯}{n}_2\underset{¯}{n}_1U)\overline{z}]/(2L_1U_2)},$$
(39)
where $`\underset{¯}{n}_i=n_i+\rho _i`$, for $`i=1,2`$. The construction of the mode functions for the orbifold $`𝒯_2/Z_N`$ is done by observing that all of the arguments in remain valid if $`n_i`$ is replaced by $`\underset{¯}{n}_i=n_i+\rho _i`$. The basis functions one is led to are given by
$`h_{\underset{¯}{n}_1,\underset{¯}{n}_2}^g(z)={\displaystyle \frac{1}{\sqrt{N}}}\eta _{\underset{¯}{n},\underset{¯}{0}}{\displaystyle \underset{k=0}{\overset{N1}{}}}g^kf_{\underset{¯}{n}_1,\underset{¯}{n}_2}(\tau ^kz)`$ (40)
where $`\eta _{\underset{¯}{n},\underset{¯}{0}}=1/\sqrt{N}`$ if $`\underset{¯}{n}_1=\underset{¯}{n}_2=0`$, and otherwise equals 1. Using this definition, one can check that the mode functions satisfy the orbifold condition
$`h_{\underset{¯}{n}_1,\underset{¯}{n}_2}^g(\tau ^kz)=g^kh_{\underset{¯}{n}_1,\underset{¯}{n}_2}^g(z).`$ (41)
As outlined in not all the functions $`h_{\underset{¯}{n}_1\underset{¯}{n}_2}`$ are independent, since they are related by
$`h_{\omega ^k(\tau )(\underset{¯}{n}_1,\underset{¯}{n}_2)}^g(z)=h_{\underset{¯}{n}_1,\underset{¯}{n}_2}^g(\tau ^kz),`$ (42)
where $`\omega ^k(\tau )(\underset{¯}{n}_1,\underset{¯}{n}_2)`$ is defined as a rotation which takes $`\underset{¯}{n}_1+\tau \underset{¯}{n}_2`$ into $`\underset{¯}{n}_1^{}+\tau \underset{¯}{n}_2^{}=\tau ^k(\underset{¯}{n}_1+\tau \underset{¯}{n}_2)`$. The set of all such rotations (i.e. for $`k=0,1,..,N1`$) identify $`N`$ domains whose union covers the whole complex plane $`(Ox;Oy)`$ defined by $`(Ox;Oy)=(\underset{¯}{n}_1+\tau _1\underset{¯}{n}_2;\tau _2\underset{¯}{n}_2)`$. Each such domain fixes the set of levels $`\underset{¯}{n}_1,\underset{¯}{n}_2`$ which identify the set of independent $`h_{\underset{¯}{n}_1,\underset{¯}{n}_2}^g`$ which are not related by the rotation $`\omega `$ of (42). This gives $`\underset{¯}{n}_1=n_1+\rho _1<0`$ and $`\underset{¯}{n}_2=n_2+\rho _20`$ as an independent set. Since here $`0\rho _{1,2}<1`$, one concludes that the conditions $`n_1<0,n_20`$ define an independent set of functions $`h_{\underset{¯}{n}_1,\underset{¯}{n}_2}`$ for the orbifold $`𝒯_2/Z_N`$.
Using the above considerations, one has the following mode decomposition for $`\mathrm{\Phi }^g`$:
$`\mathrm{\Phi }^g(x,z)={\displaystyle \underset{n_1<0,n_20}{}}\mathrm{\Phi }_{\underset{¯}{n}_1,\underset{¯}{n}_2}^g(x)h_{\underset{¯}{n}_1,\underset{¯}{n}_2}(z)+{\displaystyle \frac{\delta ^{g,1}}{\sqrt{𝒜}}}\mathrm{\Phi }_{\underset{¯}{0},\underset{¯}{0}}^1,\underset{¯}{n}_i=n_i+\rho _i,`$ (43)
which satisfies the desired condition $`\mathrm{\Phi }^g(x,\tau ^kz)=g^k\mathrm{\Phi }^g(x,z)`$. For example, for $`𝒯_2/Z_3`$
$`g=\tau ^0,\mathrm{\Phi }^0(x,z)`$ $`=`$ $`{\displaystyle \underset{n_1<0,n_20}{}}\varphi _{\underset{¯}{n}_1,\underset{¯}{n}_2}^0(x)\left[{\displaystyle \frac{1}{\sqrt{3}}}{\displaystyle \underset{k=0}{\overset{2}{}}}f_{\underset{¯}{n}_1,\underset{¯}{n}_2}(\tau ^kz)\right]+{\displaystyle \frac{1}{\sqrt{𝒜}}}\varphi _{\underset{¯}{0},\underset{¯}{0}}^0(x)`$
$`g=\tau ^1,\mathrm{\Phi }^1(x,z)`$ $`=`$ $`{\displaystyle \underset{n_1<0,n_20}{}}\varphi _{\underset{¯}{n}_1,\underset{¯}{n}_2}^1(x)\left[{\displaystyle \frac{1}{\sqrt{3}}}{\displaystyle \underset{k=0}{\overset{2}{}}}\tau ^kf_{\underset{¯}{n}_1,\underset{¯}{n}_2}(\tau ^kz)\right]`$
$`g=\tau ^2,\mathrm{\Phi }^2(x,z)`$ $`=`$ $`{\displaystyle \underset{n_1<0,n_20}{}}\varphi _{\underset{¯}{n}_1,\underset{¯}{n}_2}^2(x)\left[{\displaystyle \frac{1}{\sqrt{3}}}{\displaystyle \underset{k=0}{\overset{2}{}}}\tau ^{2k}f_{\underset{¯}{n}_1,\underset{¯}{n}_2}(\tau ^kz)\right]`$ (44)
The main difference from the $`𝒯_2/Z_2`$ orbifold is that for $`𝒯_2/Z_N`$ the sum over the Kaluza-Klein levels is restricted to positive/negative values of $`n_{1,2}`$, unlike in $`𝒯_2/Z_2`$ where one sum could be extended to the whole set $`𝐙`$ of integers. The above mode expansion leads to the following expression for the Casimir energy of a complex 6D scalar field on $`𝒯_2/Z_N`$, $`N>2`$.
$`V_𝒪(\rho _1,\rho _2)`$ $`=`$ $`{\displaystyle \underset{n_1<0,n_20}{}}\mu ^{4d}{\displaystyle \frac{d^dp}{(2\pi )^d}\mathrm{ln}\left[\frac{p^2+M_{n_1,n_2}^2(\rho _1,\rho _2)+m^2}{\mu ^2}\right]}`$ (45)
where for each $`g`$ one uses the values of $`\rho _{1,2}`$ which respect the consistency conditions.
The above domain of summation for $`n_{1,2}`$ makes the analytical calculation of $`V_𝒪(\rho _1,\rho _2)`$ more difficult than in the case of $`𝒯_2/Z_2`$. The difficulty is caused by the fact that none of the sums over $`n_1`$, $`n_2`$ can be extended<sup>4</sup><sup>4</sup>4with some exceptions in the case of $`𝒯_2/Z_4`$ orbifolds, see later. to a sum over the whole set $`𝐙`$ of integers (as we had for $`𝒯_2/Z_2`$, eq.(3.1)). As a result no (Poisson) resummation of individual contributions to $`V(\rho _1,\rho _2)`$ is possible and the calculation is then more tedious. Although one may still be able to work on the covering torus<sup>5</sup><sup>5</sup>5For a general approach to computing traces on orbifold spaces see . rather than in the orbifold basis, the approach below (being valid for any $`\rho _{1,2}`$) allows a simultaneous analysis of all orbifolds $`𝒯_2/Z_N`$, $`N>2`$.
After a long calculation (see Appendix C, eqs.(C-1) to (C-9)) one has for $`V(\rho _1,\rho _2)`$ of (45) the following result (which is valid for $`m1/L`$)
$`V_𝒪(\rho _1,\rho _2)`$ $`=`$ $`{\displaystyle \frac{1}{2𝒜^2}}\{\stackrel{~}{𝒟}+{\displaystyle \frac{m^6𝒜^3}{768\pi ^3}}{\displaystyle \frac{2}{ϵ}}+{\displaystyle \frac{(2\pi U_2)^3}{180}}[{\displaystyle \frac{1}{21}}\rho _1^2(15\rho _1^22\rho _1^4+6\rho _1^3)]`$
$`+`$ $`{\displaystyle \underset{n_1<0}{}}[(n_1+\rho _1)^2\text{Li}_3(\sigma _{n_1})+{\displaystyle \frac{3|n_1+\rho _1|}{2\pi U_2}}\text{Li}_4(\sigma _{n_1})+{\displaystyle \frac{3}{4\pi ^2U_2^2}}\text{Li}_5(\sigma _{n_1})+c.c.]\}`$ (46)
where
$`\sigma _{n_1<0}=e^{2i\pi \left(U(n_1+\rho _1)\rho _2\right)},0\rho _i<1,U=U_1+iU_2,U_2>0.`$ (47)
This is the result for the Casimir energy for $`𝒯_2/Z_N`$, $`N>2`$ with boundary conditions as in (3.2) and with $`\rho _{1,2}`$ taking the values required by the consistency conditions specific to each orbifold. Finally, $`\stackrel{~}{𝒟}`$ of (3.2) is an asymptotic series given by (see Appendix C, eq.(C-10))
$$\stackrel{~}{𝒟}=\frac{(\mu ^2𝒜)^{\frac{ϵ}{2}}}{\pi ^{ϵ/22}}\underset{n_1<0,p0}{}\frac{2(1)^p}{p!U_2^{2p2}}\mathrm{\Gamma }\left[p2+\frac{ϵ}{2}\right]\zeta [2p,\rho _2U_1(n_1+\rho _1)]\left[\frac{m^2𝒜}{(2\pi )^2U_2}+(n_1+\rho _1)^2\right]^{2p\frac{ϵ}{2}}$$
(48)
where $`\zeta [q,x]`$ is the Hurwitz zeta function , $`𝒜=L^2\mathrm{sin}\theta `$, $`U_2=\mathrm{Im}U=\mathrm{sin}\theta `$, $`\theta =2\pi /N`$. This expression of $`\stackrel{~}{𝒟}`$ is valid without any restrictions on the relative values of $`m`$, $`L`$ or $`U`$.
The quantity $`\stackrel{~}{𝒟}`$ is of particular interest because it contains additional poles as $`ϵ0`$, and so potentially introduces new contributions to the UV divergent part of $`V_𝒪`$ in (3.2). Note that if one of the sums (say that over $`n_20`$) in $`V_𝒪`$ of (45) were extended to the whole $`𝐙`$ set of integers, the quantity $`\stackrel{~}{𝒟}`$ given above would not arise due to the cancellation against the similar contribution to $`V_𝒪`$, coming from $`n_2<0`$. The latter would actually be equal to $`\stackrel{~}{𝒟}`$ of (48) with the substitutions $`\rho _21\rho _2`$ and $`U_1U_1`$. The sum of these two contributions would then vanish
$`\stackrel{~}{𝒟}(\rho _1,\rho _2)+\stackrel{~}{𝒟}(\rho _1,1\rho _2)|_{U_1U_1}=0`$ (49)
since $`\zeta [2m,x]+\zeta [2m,1x]=0`$. In particular, this explains the absence of $`\stackrel{~}{𝒟}`$ in $`𝒯_2`$ and $`𝒯_2/Z_2`$ where one of the KK sums was over the whole set $`𝐙`$. To conclude, the presence of $`\stackrel{~}{𝒟}`$ in the potential $`V_𝒪`$ is due to the fact that both sums over the Kaluza-Klein modes in (45) were restricted to positive/negative modes only.
#### 3.2.1 Ultraviolet Divergences
Because the contribution $`\stackrel{~}{𝒟}`$ potentially introduces new UV divergences for $`𝒯_2/Z_N`$ orbifolds with $`N>2`$, in this section we investigate their form in more detail in order to see what kinds of counterterms they require. In particular, we show that for $`𝒯_2/Z_{2N}`$ no new counterterms are required beyond those which already arise for the $`Z_2`$ orbifold. To do so we consider in detail the case of $`𝒯_2/Z_4`$, for which the analysis of $`\stackrel{~}{𝒟}`$ is considerably simplified. For the remaining cases $`𝒯_2/Z_N`$ (with $`N=3,6`$) the analysis follows the same technical steps as below, but is more involved and will be presented elsewhere .
In the case of $`𝒯_2/Z_4`$ which has $`U_1=0`$, $`U_2=1`$, the Hurwitz zeta function in (48) has no dependence on $`n_1`$, and this simplifies the identification of the additional poles. In the last bracket in eq.(48) one can then use a binomial expansion ($`mL1`$) or an asymptotic expansion ($`mL1`$) and following the technical details in Appendix C, eqs.(C-11), (C-12), (C) one obtains from (48) that, for $`𝒯_2/Z_4`$
$`\stackrel{~}{𝒟}`$ $`=`$ $`{\displaystyle \frac{𝒜^2}{8\pi ^2}}\{{\displaystyle \frac{m^4}{ϵ}}c_1+{\displaystyle \frac{2m^2}{ϵ}}{\displaystyle \frac{(2\pi )^2}{𝒜}}c_2+{\displaystyle \frac{1}{ϵ}}{\displaystyle \frac{(2\pi )^4}{𝒜^2}}c_3+\stackrel{~}{𝒟}_f+𝒪(ϵ)\}`$
$`c_3`$ $`=`$ $`\zeta [4,\rho _2](1/2\rho _1)+\zeta [4,\rho _1](1/2\rho _2)+2\zeta [2,\rho _1]\zeta [2,\rho _2],`$
$`c_2`$ $`=`$ $`(1/2\rho _1)\zeta [2,\rho _2]+\zeta [2,\rho _1](1/2\rho _2),c_1=(1/2\rho _1)(1/2\rho _2)`$ (50)
In the first expression $`\stackrel{~}{𝒟}_f`$ describes the terms $`𝒪(ϵ^0)`$, and its exact value depends on whether $`mL`$ is smaller or larger than unity, and is discussed later on. The zeta functions appearing in the coefficients $`c_i`$ are given by
$`\zeta [2,x]`$ $`=`$ $`{\displaystyle \frac{1}{6}}x(x1)(2x1)`$
$`\zeta [4,x]`$ $`=`$ $`{\displaystyle \frac{1}{30}}x(x1)(2x1)(3x^23x1)`$ (51)
Writing
$`V_𝒪=V_𝒪\mathrm{}+V_𝒪^{\mathrm{ren}}`$ (52)
we therefore identify the following UV divergent terms:
$`V_𝒪\mathrm{}={\displaystyle \frac{m^6𝒜}{768\pi ^3ϵ}}+{\displaystyle \frac{m^4c_1}{16\pi ^2ϵ}}+{\displaystyle \frac{m^2c_2}{2𝒜ϵ}}+{\displaystyle \frac{\pi ^2c_3}{𝒜^2ϵ}}`$ (53)
Notice that the structure of these divergences is valid independent of the relative size of $`m`$ and $`𝒜`$. For the $`Z_4`$ orbifold we have seen that consistency of the boundary conditions requires we choose the value of $`\rho _1=\rho _2=\rho `$ with $`\rho `$ equal to 0 or $`\frac{1}{2}`$, and so we must evaluate the coefficients $`c_k`$ with these choices. Since both $`\zeta [2,\rho ]`$ and $`\zeta [4,\rho ]`$ vanish when $`\rho =0`$ or $`\rho =\frac{1}{2}`$, we see that for all such cases
$`c_2=c_3=0`$ (54)
leaving in $`V_𝒪\mathrm{}`$ only the divergences in the first and second terms in eq.(53).
The first term in (53) is a renormalization of the bulk cosmological constant and is present irrespective of the values of $`\rho _1`$ or $`\rho _2`$. Its coefficient is $`1/4`$ the size of the similar result for the covering torus. Therefore, once this term is expressed in terms of the orbifold area, $`𝒜_𝒪=𝒜/4`$, its coefficient is precisely the same as was found for $`𝒯_2`$ and $`𝒯_2/Z_2`$, as expected. The term in (53), proportional to $`c_1`$, is a “brane” divergence, which renormalizes the brane tension. It is nonzero only for the case where $`\rho _1=\rho _2=0`$, in which case $`c_1=\frac{1}{4}`$.
To conclude, we find that the UV divergences for the Casimir energy due to compactifications on $`𝒯_2/Z_4`$ with discrete Wilson lines $`(\rho _1,\rho _2)`$, have two kinds of divergences at one loop, similar to the case of $`𝒯_2/Z_2`$. These have the form and coefficients required by the general heat-kernel analysis and renormalize the bulk cosmological constant and the tension of branes localized at the orbifold fixed points. For the case of remaining orbifolds $`𝒯_2/Z_3`$, $`𝒯_2/Z_6`$, the analysis of the divergences of the quantity $`\stackrel{~}{𝒟}`$ of (48) is more involved since the Zeta function entering its definition will retain a $`n_1`$ dependence. This makes the computation more tedious and the identification of the relevant counterterms more difficult to analyze in this case .
#### 3.2.2 The finite part of Casimir Energy.
For the finite part of the Casimir energy for the orbifold $`𝒯_2/Z_4`$ one obtains in the limit of $`m^2𝒜1`$ (see Appendix eqs.(C-12) and (C)).
$`V_𝒪^{\mathrm{ren}}(\rho _1,\rho _2)`$ $`=`$ $`{\displaystyle \frac{1}{2𝒜^2}}\{\stackrel{~}{𝒟}_f+{\displaystyle \frac{2\pi ^3}{45}}[{\displaystyle \frac{1}{21}}\rho _1^2(15\rho _1^22\rho _1^4+6\rho _1^3)]`$
$`+`$ $`{\displaystyle \underset{n_1<0}{}}[(n_1+\rho _1)^2\text{Li}_3(\sigma _{n_1})+{\displaystyle \frac{3|n_1+\rho _1|}{2\pi }}\text{Li}_4(\sigma _{n_1})+{\displaystyle \frac{3}{4\pi ^2}}\text{Li}_5(\sigma _{n_1})+c.c.]\}`$
with the notation
$`\sigma _{n_1<0}=e^{2\pi \left(n_1+\rho _1+i\rho _2\right)},0\rho _i<1`$ (56)
Here $`\stackrel{~}{𝒟}_f`$ is an asymptotic series which for $`m^2𝒜1`$ has the following expression (see Appendix C, eq.(C-12))
$`\stackrel{~}{𝒟}_f`$ $`=`$ $`\pi ^2\{2\zeta [2,\rho _1]\zeta [2,\rho _2]\mathrm{ln}(\tau \pi e^{\gamma 1})+\zeta [0,\rho _2](\zeta [4,\rho _1]\mathrm{ln}(\pi \tau e^{\gamma \frac{3}{2}})+2\zeta ^{}[4,1\rho _1])`$ (57)
$`+`$ $`4\zeta [2,\rho _2]\zeta ^{}[2,1\rho _1]+\zeta [4,\rho _2]\left(\zeta [0,\rho _1]\mathrm{ln}(\pi \tau e^\gamma )+2\zeta ^{}[0,1\rho _1]\right)`$
$`+`$ $`2{\displaystyle \underset{k3}{}}{\displaystyle \frac{(1)^k}{k!}}\mathrm{\Gamma }[k2]\zeta [2k,\rho _2]\zeta [2k4,1\rho _1]\},\mathrm{with}\tau =(2\pi )^2/(\mu ^2𝒜).`$
Since<sup>6</sup><sup>6</sup>6In (57) the derivative of Zeta function is taken wrt its first argument. for the orbifold $`𝒯_2/Z_4`$ the Wilson lines have the values $`\rho _1=\rho _2=0,1/2`$, $`\stackrel{~}{𝒟}_f`$ simplifies to give:
$`\stackrel{~}{𝒟}_f|_{\rho _1=\rho _2=0}={\displaystyle \frac{3}{4\pi ^2}}\zeta [5],\mathrm{and}\stackrel{~}{𝒟}_f|_{\rho _1=\rho _2=\frac{1}{2}}=0.`$ (58)
Eqs.(3.2.2), (58) and also (52), (53), (54) give the final result for the Casimir energy for the orbifold $`𝒯_2/Z_4`$ with discrete Wilson lines. Finally, with $`\rho _1=\rho _2=0,1/2`$ one has from eqs.(3.2.2) to (58)
$`V_𝒪^{\mathrm{ren}}(\rho ,\rho )={\displaystyle \frac{1}{4}}V^{\mathrm{ren}}(\rho ,\rho )|_{m=0},\rho =0,{\displaystyle \frac{1}{2}}`$ (59)
where $`V^{\mathrm{ren}}`$ is the result for the 2-torus $`𝒯_2`$, given in eq.(2.1.2). Following closely these steps, one can also obtain from eqs.(3.2), (48) similar results for $`𝒯_2/Z_6`$ and $`𝒯_2/Z_3`$ orbifolds.
#### 3.2.3 Heavy-mass dependence.
We discuss now the heavy mass dependence for the Casimir energy. For $`𝒯_2/Z_4`$ with $`m1/L`$ it turns out that the divergences in $`V_𝒪`$ are identical to those in eq.(53). From Appendix C, eq.(C-15) with (C-9), (C), (C-14) one obtains the full result for $`V_𝒪`$ for $`mL1`$. Here we outline only the main behaviour which is
$`V_𝒪`$ $`=`$ $`\left\{{\displaystyle \frac{m^4c_1}{32\pi ^2}}\mathrm{ln}(\pi e^{\gamma \frac{3}{2}}m^2/\mu ^2)+{\displaystyle \frac{1}{2}}{\displaystyle \frac{m^6𝒜}{768\pi ^3}}\mathrm{ln}(\pi e^{\gamma \frac{11}{6}}m^2/\mu ^2)+\mathrm{}\right\}`$ (60)
with $`c_1=(1/2\rho )^2`$ and where the dots account for additional terms such as polylogarithms terms, identical to those in (3.2.2), and for (asymptotic series of) terms which are suppressed by inverse powers of $`m^2𝒜`$. The latter vanish in the special case of $`𝒯_2/Z_4`$ with $`\rho =0,1/2`$, to leave only the (exponentially suppressed) polylogarithm contributions.
## 4 Conclusions
In this paper we compute the value of the Casimir energy for a very broad class of two-dimensional toroidal compactifications. These include the general case of $`𝒯_2`$ compactifications with arbitrary boundary conditions for the 6D fields corresponding to the presence of arbitrary Wilson lines, as well as $`𝒯_2/Z_N`$ orbifolds (also with Wilson lines) obtained by identifying points under $`Z_N`$ rotations. Our calculations are explicit for a 6D scalar having an arbitrary 6D mass $`m`$, and we show how to extend these results to higher-spin fields for supersymmetric 6D theories. Particular attention was paid to regularization issues and to the identification of the divergences of the potential. The computation also investigated the dependence of the result on $`m`$, including limits for which $`m^2𝒜`$ is larger or smaller than unity, (where $`𝒜`$ is the volume of the internal 2 dimensions).
For the cases of $`𝒯_2`$ and $`𝒯_2/Z_2`$, our calculation generalizes earlier results to include the dependence on an arbitrary complex structure, $`U`$, for the underlying torus. The potential obtained is likely to be useful for studies of the dynamics of these moduli, including their stabilization and their potential applications to cosmology .
By carefully isolating the UV divergent part of $`V`$, we show that all of the divergences may be renormalized into a bulk cosmological constant (which gives a Casimir energy proportional to $`m^6𝒜`$) and - for the case of $`𝒯_2/Z_2`$ \- a cosmological constant (or brane tension) localized at the orbifold fixed points (which gives a Casimir energy proportional to $`m^4`$). Furthermore, these divergences agree with expectations based on general heat-kernel calculations, such as those recently performed for 6D compactifications in ref. . For massive 6D scalar fields, $`m^2𝒜1`$, the dependence on $`m`$ of the finite part of the Casimir energy obtained in the modified minimal subtraction scheme, is exponentially suppressed.
We present results for the Casimir energy also for $`𝒯_2/Z_N`$ orbifolds with $`N>2`$, again including Wilson lines and any shape moduli which are allowed. The case $`𝒯_2/Z_4`$ was studied in particular detail. The UV divergences that emerge in this case again take the form required by general heat-kernel arguments, and can be absorbed into renormalizations of the bulk cosmological constant and brane tensions localized at the orbifold fixed points. The finite part of the Casimir energy was computed in detail and may be used for phenomenological applications. Finally, the technical tools of the Appendix can be used for other applications such as the one-loop corrections to the gauge couplings in gauge theories on orbifolds, in the presence of discrete Wilson lines.
Acknowledgements.
We thank Y. Aghababaie, Z. Chacko, J. Elliot, G. Gabadadze and A. de la Macorra for helpful discussions about 6D Casimir energies on the torus. D.G. thanks S. Groot-Nibbelink and Hyun Min Lee for discussions on some technical aspects of this work.
C.B.’s research is supported by grants from NSERC (Canada) and McMaster University and D.H. acknowledges partial support from McGill University. The research of F.Q. is partially supported by PPARC and a Royal Society Wolfson award. The work of D. Ghilencea was supported by a post-doctoral research grant from Particle Physics and Astronomy Research Council (PPARC), United Kingdom. D. Ghilencea acknowledges the support from the RTN European program MRTN-CT-2004-503369, to attend the “Planck 2005” conference where this work was completed.
## Appendix
### A . Calculation of the vacuum energy in DR for 2D compactifications.
We provide here details of the calculation of the vacuum energy. One has ($`d=4ϵ`$)
$`V^{}(\rho _1,\rho _2)\mu ^{4d}{\displaystyle \underset{n_{1,2}𝐙}{\overset{}{}}}{\displaystyle \frac{d^dp}{(2\pi )^d}\mathrm{ln}\left[p^2+M_{n_1,n_2}^2\right]}={\displaystyle \frac{\mu ^4}{(2\pi )^d}}{\displaystyle \underset{n_{1,2}𝐙}{\overset{}{}}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t^{1+d/2}}}e^{\pi t\left[M_{n_1,n_2}^2/\mu ^2\right]}`$ (A-1)
$`\mu `$ is a finite, non-zero mass scale introduced by the DR scheme. A “prime” on a double sum excludes the $`(n_1,n_2)=(0,0)`$ mode. If a level $`(n_1,n_2)`$ is massless $`M_{n_1,n_2}=0`$ (for example if $`\rho _{1,2}𝐙`$)), mathematical consistency requires one shift $`M_{n_1,n_2}^2M_{n_1,n_2}^2+m^2`$ by a finite non-zero $`m^2=\delta \mu ^2`$ ($`\delta `$ dimensionless). This also helps us identify the scale ($`m`$) dependence of the divergences (poles in $`ϵ`$). We use
$$M_{n_1,n_2}^2=\frac{(2\pi )^2}{𝒜U_2}|n_2+\rho _2U(n_1+\rho _1)|^2;UU_1+iU_2=e^{i\theta }\frac{L_2}{L_1};𝒜=L_1L_2\mathrm{sin}\theta $$
(A-2)
The DR regularized sum in (A-1) is re-written
$`V^{}(\rho _1,\rho _2)=\mu ^4C_ϵ{\displaystyle \underset{n_{1,2}𝐙}{\overset{}{}}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t^{3ϵ/2}}}e^{\pi t[M_{n_1,n_2}^2/\mu ^2+\delta ]}\mu ^4C_ϵ𝒥_{ϵ/2}^{},C_ϵ={\displaystyle \frac{1}{(2\pi )^{4ϵ}}}`$ (A-3)
with $`\delta =m^2/\mu ^2`$. The calculation of $`V^{}`$ is reduced to that of $`𝒥_ϵ^{}`$ performed below.
### B . Series of Kaluza-Klein integrals and their DR regularization.
We evaluate (in the text $`\tau (2\pi )^2/(\mu ^2𝒜U_2)`$, $`\delta m^2/\mu ^2`$)
$`𝒥_ϵ^{}{\displaystyle \underset{n_1,n_2𝐙}{\overset{}{}}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t^{3+ϵ}}}e^{\pi t\tau |n_2+\rho _2U(n_1+\rho _1)|^2}e^{\pi \delta t},\tau >0,UU_1+iU_2,\delta >0,\rho _i𝐑`$ (B-1)
$`𝒥_ϵ^{}`$ includes shape moduli effects ($`\theta \pi /2`$, $`L_1L_2`$), $`\delta `$ shifts, and arbitrary “twists” $`\rho _i`$ wrt $`L_{1,2}`$. To evaluate $`𝒥_ϵ^{}`$ one uses re-summation (E-1); the integrand becomes
$`{\displaystyle \underset{n_1,n_2𝐙}{\overset{}{}}}e^{\pi t\tau |n_2+\rho _2U(n_1+\rho _1)|^2}`$ $`=`$ $`{\displaystyle \underset{n_2𝐙}{\overset{}{}}}e^{\pi t\tau |n_2+\rho _2U\rho _1|^2}+{\displaystyle \underset{n_1𝐙}{\overset{}{}}}{\displaystyle \underset{n_2𝐙}{}}e^{\pi t\tau |n_2+\rho _2U(n_1+\rho _1)|^2}`$
$`=`$ $`{\displaystyle \underset{n_2𝐙}{\overset{}{}}}e^{\pi t\tau |n_2+\rho _2U\rho _1|^2}+{\displaystyle \frac{1}{\sqrt{t\tau }}}{\displaystyle \underset{n_1𝐙}{\overset{}{}}}e^{\pi t\tau U_2^2(n_1+\rho _1)^2}`$
$`+`$ $`{\displaystyle \frac{1}{\sqrt{t\tau }}}{\displaystyle \underset{n_1𝐙}{\overset{}{}}}{\displaystyle \underset{\stackrel{~}{n}_2𝐙}{\overset{}{}}}e^{\frac{\pi \stackrel{~}{n}_{2}^{}{}_{}{}^{2}}{t\tau }\pi t\tau U_2^2(n_1+\rho _1)^2+2\pi i\stackrel{~}{n}_2(\rho _2U_1(\rho _1+n_1))}`$ (B-2)
A prime on a double sum indicates that $`(n_1,n_2)(0,0)`$ is excluded and a “prime” on a single sum excludes its $`n=0`$ mode. The three contributions above can be integrated termwise for any real $`\rho _i`$, (given the presence of $`e^{\pi t\delta }`$). Accordingly, one has three contributions
$`𝒥_ϵ^{}(\rho _1,\rho _2)`$ $`=`$ $`𝒦_1(\rho _1,\rho _2)+𝒦_2(\rho _1,\rho _2)+𝒦_3(\rho _1,\rho _2)`$ (B-3)
defined/evaluated in the following:
$``$ Computing $`𝒦_1`$:
$`𝒦_1`$ $``$ $`{\displaystyle \underset{n_2𝐙}{\overset{}{}}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t^{3+ϵ}}}e^{\pi t\tau |n_2+\rho _2U\rho _1|^2}e^{\pi \delta t}={\displaystyle \frac{\pi ^2}{2ϵ}}(\delta +\tau |\rho _2U\rho _1|^2)^2`$
$`+`$ $`{\displaystyle \frac{\pi ^2}{4}}(\delta +\tau |\rho _2U\rho _1|^2)^2\mathrm{ln}\left[\pi ^2e^{2\gamma 3}(\delta +\tau |\rho _2U\rho _1|^2)^2\right]{\displaystyle \frac{8\pi ^3}{15\sqrt{\tau }}}(\delta +\tau U_2^2\rho _1^2)^{\frac{5}{2}}`$
$`+`$ $`\tau (\delta +\tau U_2^2\rho _1^2)\text{Li}_3(e^{2\pi \gamma (0)})+{\displaystyle \frac{3\tau ^{\frac{3}{2}}}{2\pi }}(\delta +\tau U_2^2\rho _1^2)^{\frac{1}{2}}\text{Li}_4(e^{2\pi \gamma (0)})+{\displaystyle \frac{3\tau ^2}{4\pi ^2}}\text{Li}_5(e^{2\pi \gamma (0)})+c.c.`$
where “c.c.” applies to the PolyLogarithm functions only. To evaluate $`𝒦_1`$ we first added and subtracted the $`n_2=0`$ mode contribution. We then used a (Poisson) re-summation over $`n_2`$, then the integral representation of modified Bessel functions (E-2) with (E), and the definition of the Polylogarithm $`\text{Li}_\sigma (x)`$ (E-4). Finally we used the notation
$`\gamma (0){\displaystyle \frac{1}{\sqrt{\tau }}}(\delta +\tau U_2^2\rho _1^2)^{\frac{1}{2}}i(\rho _2U_1\rho _1)`$ (B-5)
The divergence of $`𝒦_1`$ is that of the excluded $`n_2=0`$ mode in $`𝒦_1`$, which is in turn due to the absence of $`(n_1,n_2)=(0,0)`$ in the definition of $`𝒥_ϵ^{}`$.
$``$ Computing $`𝒦_2`$: We introduce the notation $`\mathrm{\Delta }_{\rho _1}\rho _1[\rho _1]`$, $`0\mathrm{\Delta }_{\rho _1}<1`$, $`[\rho _1]𝐙`$.
(a). For $`0\delta /(\tau U_2^2)<1`$ we have
$`𝒦_2`$ $``$ $`{\displaystyle \frac{1}{\sqrt{\tau }}}{\displaystyle \underset{n_1𝐙}{\overset{}{}}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t^{7/2+ϵ}}}e^{\pi t\tau U_2^2(n_1+\rho _1)^2}e^{\pi \delta t}`$
$`=`$ $`{\displaystyle \frac{1}{\sqrt{\tau }}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t^{7/2+ϵ}}}\left[{\displaystyle \underset{n_1𝐙}{}}e^{\pi t[\delta +\tau U_2^2(n_1+\mathrm{\Delta }_{\rho _1})^2]}e^{\pi t(\delta +\tau U_2^2\rho _1^2)}\right]`$
$`=`$ $`{\displaystyle \frac{\pi ^{5/2+ϵ}}{\sqrt{\tau }}}\mathrm{\Gamma }[5/2ϵ]\left[{\displaystyle \underset{n_1𝐙}{\overset{}{}}}[\delta +\tau U_2^2(n_1+\mathrm{\Delta }_{\rho _1})^2]^{\frac{5}{2}+ϵ}(\delta +\tau U_2^2\rho _1^2)^{\frac{5}{2}+ϵ}+(\delta +\tau U_2^2\mathrm{\Delta }_{\rho _1}^2)^{\frac{5}{2}+ϵ}\right]`$
$`=`$ $`{\displaystyle \frac{\pi ^{\frac{5}{2}}}{\sqrt{\tau }}}\mathrm{\Gamma }[5/2]\left[(\delta +\tau U_2^2\mathrm{\Delta }_{\rho _1}^2)^{\frac{5}{2}}(\delta +\tau U_2^2\rho _{1}^{}{}_{}{}^{2})^{\frac{5}{2}}\right]`$
$`+`$ $`{\displaystyle \frac{(\pi \tau U_2^2)^{\frac{5}{2}+ϵ}}{\sqrt{\tau }}}{\displaystyle \underset{k0}{}}{\displaystyle \frac{\mathrm{\Gamma }[k5/2ϵ]}{k!}}\left[{\displaystyle \frac{\delta }{\tau U_2^2}}\right]^k\left[\zeta [2k52ϵ,1+\mathrm{\Delta }_{\rho _1}]+(\mathrm{\Delta }_{\rho _1}\mathrm{\Delta }_{\rho _1})\right]`$ (B-6)
In the last step we used the binomial expansion
$$\underset{n0}{}[a(n+c)^2+q]^s=a^s\underset{k0}{}\frac{\mathrm{\Gamma }[k+s]}{k!\mathrm{\Gamma }[s]}\left[\frac{q}{a}\right]^k\zeta [2k+2s,c],0<q/a1$$
(B-7)
Here $`\zeta [q,a]`$ with $`a0,1,2,\mathrm{}`$ is the Hurwitz zeta function, (with $`\zeta [q,a]=_{n0}(a+n)^q`$ for $`\text{Re}(q)>1`$). Hurwitz zeta-function has one singularity (simple pole) at $`q=1`$ and $`\zeta [q,1]=\zeta [q]`$ with $`\zeta [q]`$ the Riemann zeta function. The only divergence in $`𝒦_2`$ is due to its $`k=3`$ term in (B-6), from the singularity of the Zeta function. In the remaining terms in the series one can safely set $`ϵ=0`$. Further
$`\zeta [12ϵ,1\pm \mathrm{\Delta }_{\rho _1}]`$ $`=`$ $`{\displaystyle \frac{1}{2ϵ}}\psi (1\pm \mathrm{\Delta }_{\rho _1})+𝒪(ϵ)`$
$`\mathrm{\Gamma }[1/2ϵ]`$ $`=`$ $`\pi ^{1/2}(1+ϵ\mathrm{ln}(4e^\gamma ))+𝒪(ϵ)`$
$`x^ϵ`$ $`=`$ $`1+ϵ\mathrm{ln}x+𝒪(ϵ)`$ (B-8)
we find for $`0\delta /(\tau U_2^2)<1`$:
$`𝒦_2`$ $`=`$ $`{\displaystyle \frac{\pi ^3\delta ^3}{6\tau |U_2|}}{\displaystyle \frac{1}{ϵ}}+{\displaystyle \frac{\pi ^3\delta ^3}{6\tau |U_2|}}\mathrm{ln}\left[4\pi \tau U_2^2e^{\gamma +\psi (\mathrm{\Delta }_{\rho _1})+\psi (\mathrm{\Delta }_{\rho _1})}\right]`$
$``$ $`{\displaystyle \frac{8\pi ^3}{15\sqrt{\tau }}}\left[(\delta +\tau U_2^2\mathrm{\Delta }_{\rho _1}^2)^{\frac{5}{2}}(\delta +\tau U_2^2\rho _{1}^{}{}_{}{}^{2})^{\frac{5}{2}}+(\tau U_2^2)^{\frac{5}{2}}(\zeta [5,1+\mathrm{\Delta }_{\rho _1}]+\zeta [5,1\mathrm{\Delta }_{\rho _1}])\right]`$
$`+`$ $`{\displaystyle \frac{4\pi ^3}{3\sqrt{\tau }}}\delta \tau ^{3/2}|U_2|^3(\zeta [3,1+\mathrm{\Delta }_{\rho _1}]+\zeta [3,1\mathrm{\Delta }_{\rho _1}])+\pi ^3\delta ^2|U_2|(1/6+\mathrm{\Delta }_{\rho _1}^2)`$
$`+`$ $`{\displaystyle \frac{\pi ^{5/2}}{\sqrt{\tau }}}(\tau U_2^2)^{5/2}{\displaystyle \underset{p1}{}}{\displaystyle \frac{\mathrm{\Gamma }[p+1/2]}{(p+3)!}}\left[{\displaystyle \frac{\delta }{\tau U_2^2}}\right]^{p+3}\left[\zeta [2p+1,1+\mathrm{\Delta }_{\rho _1}]+\zeta [2p+1,1\mathrm{\Delta }_{\rho _1}]\right]`$ (B-9)
The divergence in $`𝒦_2`$ is due to $`\stackrel{~}{n}_2=0`$ (i.e. Poisson re-summed zero mode wrt to the second dimension) in the presence of infinitely many KK modes of the first dimension ($`n_1`$). It is thus an interplay effect of both compact dimensions. The condition of validity of the above result $`0\delta /(\tau U_2^2)<1`$ gives $`\delta \mu ^21/L_{1,2}^2`$. If $`\delta 1`$ the result (B) simplifies considerably.
(b). If $`\delta /(\tau U_2^2)1`$ or $`\delta 1`$ eq.(B-7) does not converge. If so, $`𝒦_2`$ is reevaluated as below:
$`𝒦_2`$ $``$ $`{\displaystyle \frac{1}{\sqrt{\tau }}}{\displaystyle \underset{n_1𝐙}{\overset{}{}}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t^{7/2+ϵ}}}e^{\pi t\tau U_2^2(n_1+\rho _1)^2}e^{\pi \delta t}`$
$`=`$ $`{\displaystyle \frac{1}{\sqrt{\tau }}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t^{7/2+ϵ}}}\left[{\displaystyle \underset{n_1𝐙}{}}e^{\pi t[\delta +\tau U_2^2(n_1+\rho _1)^2]}e^{\pi t(\delta +\tau U_2^2\rho _1^2)}\right]`$
$`=`$ $`{\displaystyle \frac{1}{\sqrt{\tau }}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t^{7/2+ϵ}}}[{\displaystyle \frac{1}{\sqrt{t\tau }U_2}}{\displaystyle \underset{\stackrel{~}{n}_1𝐙}{\overset{}{}}}e^{\frac{\pi \stackrel{~}{n}_1^2}{t\tau U_2^2}+2i\pi \stackrel{~}{n}_1\rho _1\pi \delta t}e^{\pi t(\delta +\tau U_2^2\rho _1^2)}]+{\displaystyle \frac{(\pi \delta )^{3+ϵ}}{\tau U_2}}\mathrm{\Gamma }[3ϵ]`$ (B-10)
where the last term originates in the $`\stackrel{~}{n}_1=0`$ term of the series. We find (with (E-2))
$`𝒦_2`$ $`=`$ $`{\displaystyle \frac{\pi ^3\delta ^3}{6\tau |U_2|}}\left[{\displaystyle \frac{1}{ϵ}}+\mathrm{ln}\left(\pi \delta e^{\gamma 11/6}\right)\right]+4\sqrt{\tau }U_2^2\delta ^{\frac{3}{2}}{\displaystyle \underset{\stackrel{~}{n}_11}{}}{\displaystyle \frac{\mathrm{cos}[2\pi \stackrel{~}{n}_1\rho _1]}{\stackrel{~}{n}_1^3}}K_3(2\pi \stackrel{~}{n}_1\sqrt{\delta /(\tau U_2^2)})`$ (B-11)
$`+`$ $`{\displaystyle \frac{8\pi ^3}{15\sqrt{\tau }}}(\delta +\tau U_2^2\rho _1^2)^{\frac{5}{2}}`$
rapidly convergent if $`\delta \tau U_2^2`$ or $`\delta 1`$. This ends our calculation of $`𝒦_2`$ at large/small $`\delta `$.
$``$ Computing $`𝒦_3`$:
$`𝒦_3{\displaystyle \frac{1}{\sqrt{\tau }}}{\displaystyle \underset{n_1𝐙}{\overset{}{}}}{\displaystyle \underset{\stackrel{~}{n}_2𝐙}{\overset{}{}}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t^{7/2+ϵ}}}e^{\pi \stackrel{~}{n}_2^2/(t\tau )\pi t\tau U_2^2(n_1+\rho _1)^2+2\pi i\stackrel{~}{n}_2(\rho _2U_1(\rho _1+n_1))\pi \delta t}`$ (B-12)
Since $`𝒦_3`$ is always exponentially suppressed at $`t0`$ and at $`t\mathrm{}`$ it has no singularities. We can thus safely set $`ϵ=0`$. One finds
$`𝒦_3`$ $`=`$ $`\tau {\displaystyle \underset{n_1𝐙}{\overset{}{}}}\left[z(n_1)\text{Li}_3(e^{2\pi \gamma (n_1)})+{\displaystyle \frac{3}{2\pi }}(\tau z(n_1))^{1/2}\text{Li}_4(e^{2\pi \gamma (n_1)})+{\displaystyle \frac{3\tau }{4\pi ^2}}\text{Li}_5(e^{2\pi \gamma (n_1)})\right]+c.c.`$
$`z(n_1)`$ $``$ $`\delta +\tau U_2^2(n_1+\rho _1)^2`$
$`\gamma (n_1)`$ $``$ $`{\displaystyle \frac{1}{\sqrt{\tau }}}(\delta +\tau U_2^2(n_1+\rho _1)^2)^{1/2}i(\rho _2U_1(n_1+\rho _1))`$ (B-13)
To evaluate $`𝒦_3`$ we used the representation of Bessel functions eq.(E-2), then (E) and finally the polylogarithm definition in (E-4). This result simplifies considerably if $`\delta 1`$.
To conclude if $`0\delta /(\tau U_2^2)<1`$ we find for $`𝒥_ϵ^{}`$ (with (B-1), (B-3), (B), (B), (B))
$`𝒥_ϵ^{}`$ $``$ $`{\displaystyle \underset{n_1,n_2𝐙}{\overset{}{}}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t^{3+ϵ}}}e^{\pi t\tau |n_2+\rho _2U(n_1+\rho _1)|^2}e^{\pi \delta t}`$
$`=`$ $`{\displaystyle \frac{\pi ^2}{2ϵ}}(\delta +\tau |\rho _2U\rho _1|^2)^2+{\displaystyle \frac{\pi ^2}{4}}(\delta +\tau |\rho _2U\rho _1|^2)^2\mathrm{ln}\left[\pi ^2e^{2\gamma 3}(\delta +\tau |\rho _2U\rho _1|^2)^2\right]`$
$`+`$ $`\tau {\displaystyle \underset{n_1𝐙}{}}\left[z(n_1)\text{Li}_3(e^{2\pi \gamma (n_1)})+{\displaystyle \frac{3}{2\pi }}(\tau z(n_1))^{1/2}\text{Li}_4(e^{2\pi \gamma (n_1)})+{\displaystyle \frac{3\tau }{4\pi ^2}}\text{Li}_5(e^{2\pi \gamma (n_1)})\right]+c.c.`$
$``$ $`{\displaystyle \frac{8\pi ^3}{15\sqrt{\tau }}}(\delta +\tau U_2^2\mathrm{\Delta }_{\rho _1}^2)^{\frac{5}{2}}+{\displaystyle \frac{\pi ^3\delta ^3}{6\tau |U_2|}}\left[{\displaystyle \frac{1}{ϵ}}+\mathrm{ln}\left[4\pi \tau U_2^2e^{\gamma +\psi (\mathrm{\Delta }_{\rho _1})+\psi (\mathrm{\Delta }_{\rho _1})}\right]\right]+\pi ^3\delta ^2|U_2|\left({\displaystyle \frac{1}{6}}+\mathrm{\Delta }_{\rho _1}^2\right)`$
$`+`$ $`{\displaystyle \frac{4\pi ^3}{45}}\tau ^2|U_2|^5\left[{\displaystyle \frac{1}{21}}\mathrm{\Delta }_{\rho _1}^2(15\mathrm{\Delta }_{\rho _1}^22\mathrm{\Delta }_{\rho _1}^4)\right]{\displaystyle \frac{\pi ^3\delta \tau |U_2|^3}{45}}\left[130\mathrm{\Delta }_{\rho _1}^2(1+\mathrm{\Delta }_{\rho _1}^2)\right]`$
$`+`$ $`\pi ^{5/2}\tau ^2|U_2|^5{\displaystyle \underset{p1}{}}{\displaystyle \frac{\mathrm{\Gamma }[p+1/2]}{(p+3)!}}\left[{\displaystyle \frac{\delta }{\tau U_2^2}}\right]^{p+3}\left[\zeta [2p+1,1+\mathrm{\Delta }_{\rho _1}]+\zeta [2p+1,1\mathrm{\Delta }_{\rho _1}]\right]`$ (B-14)
This restriction $`0\delta /(\tau U_2^2)<1`$ is required for the convergence of the calculation of $`𝒦_2`$. The first line in $`𝒥_ϵ^{}`$ is due to the absence of the mode $`(0,0)`$. $`𝒥_ϵ^{}`$ is well defined even for $`\delta =0`$ if $`M_{n_1,n_2}0`$. In such case the result is obtained by redoing the above calculation with $`\delta =0`$ or more easily, from the one above by formally setting $`\delta =0`$. The above result simplifies considerably when $`\delta 1`$.
Finally, if we have $`\delta /(\tau U_2^2)1`$, from eqs.(B-1), (B-3), (B), (B-11), (B) we find
$`𝒥_ϵ^{}`$ $``$ $`{\displaystyle \underset{n_1,n_2𝐙}{\overset{}{}}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t^{3+ϵ}}}e^{\pi t\tau |n_2+\rho _2U(n_1+\rho _1)|^2}e^{\pi \delta t}`$
$`=`$ $`{\displaystyle \frac{\pi ^2}{2ϵ}}(\delta +\tau |\rho _2U\rho _1|^2)^2+{\displaystyle \frac{\pi ^2}{4}}(\delta +\tau |\rho _2U\rho _1|^2)^2\mathrm{ln}\left[\pi ^2e^{2\gamma 3}(\delta +\tau |\rho _2U\rho _1|^2)^2\right]`$
$`+`$ $`\tau {\displaystyle \underset{n_1𝐙}{}}\left[z(n_1)\text{Li}_3(e^{2\pi \gamma (n_1)})+{\displaystyle \frac{3}{2\pi }}(\tau z(n_1))^{1/2}\text{Li}_4(e^{2\pi \gamma (n_1)})+{\displaystyle \frac{3\tau }{4\pi ^2}}\text{Li}_5(e^{2\pi \gamma (n_1)})\right]+c.c.`$
$`+`$ $`{\displaystyle \frac{\pi ^3\delta ^3}{6\tau |U_2|}}\left[{\displaystyle \frac{1}{ϵ}}+\mathrm{ln}\left(\pi \delta e^{\gamma 11/6}\right)\right]+4\sqrt{\tau }U_2^2\delta ^{\frac{3}{2}}{\displaystyle \underset{\stackrel{~}{n}_11}{}}{\displaystyle \frac{\mathrm{cos}[2\pi \stackrel{~}{n}_1\rho _1]}{\stackrel{~}{n}_1^3}}K_3\left(2\pi \stackrel{~}{n}_1\sqrt{\delta /(\tau U_2^2)}\right)`$ (B-15)
and this concludes the evaluation of $`𝒥_ϵ^{}`$. Using (E), one shows that the last equation has the contributions from $`K_3`$ and from the polylogarithms suppressed if $`\delta \tau U_2^2`$ and $`\delta \tau `$, to leave the first line and the term $`\pi ^3\delta ^3/(6ϵ\tau |U_2|)`$ as its leading behaviour for this region of the parameter space.
$``$ Adding to $`𝒥_ϵ^{}`$ the effect of the mode $`(0,0)`$ the result is
$`𝒥_ϵ(\rho _1,\rho _2)𝒥_ϵ^{}(\rho _1,\rho _2)+𝒵_ϵ(\rho _1,\rho _2)`$ (B-16)
with
$`𝒵_ϵ(\rho _1,\rho _2)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t^{3+ϵ}}}e^{\pi t\tau |\rho _2U\rho _1)|^2}e^{\pi \delta t}`$ (B-17)
$`=`$ $`{\displaystyle \frac{\pi ^2}{2ϵ}}(\delta +\tau |\rho _2U\rho _1|^2)^2{\displaystyle \frac{\pi ^2}{4}}(\delta +\tau |\rho _2U\rho _1|^2)^2\mathrm{ln}\left[\pi ^2e^{2\gamma 3}(\delta +\tau |\rho _2U\rho _1|^2)^2\right]`$
$`𝒥_ϵ`$ is thus given by $`𝒥_ϵ^{}`$ without the first line in (B) or (B).
Eqs. (8), (2.1.2) in the text are then obtained from
$`V(\rho _1,\rho _2)=\mu ^4C_ϵ𝒥_{ϵ/2}(\rho _1,\rho _2),\mathrm{with}\delta \mathrm{m}^2/\mu ^2,\tau (2\pi )^2/(\mu ^2𝒜\mathrm{U}_2).`$ (B-18)
for the case $`\delta 1`$.
### C . More series of Kaluza-Klein integrals for $`T_2/Z_N`$ orbifolds.
For general orbifolds $`𝒯_2/Z_N`$ one needs to evaluate ($`UU_1+iU_2`$)
$`_s`$ $`=`$ $`{\displaystyle \underset{n_1<0,n_20}{}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t^{1s}}}e^{\pi t\tau |n_2+\rho _2U(n_1+\rho _1)|^2}e^{\pi t\delta }`$ (C-1)
$`=`$ $`{\displaystyle \underset{n_1<0,n_20}{}}\mathrm{\Gamma }[s]\pi ^s\left(\tau |n_2+\rho _2U(n_1+\rho _1)|^2+\delta \right)^s,\tau ,\delta >0,\mathrm{\hspace{0.17em}\hspace{0.17em}0}\rho _{1,2}<1.`$
Eq.(3.2) in the text is then $`V_𝒪=\mu ^4/(2\pi )^{4ϵ}_{(2\frac{ϵ}{2})}\left(\delta m^2/\mu ^2,\tau (2\pi )^2/(\mu ^2𝒜U_2)\right)`$.
To compute $`_s`$, the usual (Poisson) re-summation used in previous sections is not applicable given the restricted summation on $`n_1,n_2`$. The sum of the last line is actually a “truncated” Epstein function. To analyze $`_s`$, we follow the method in , for both non-zero $`\rho _{1,2}`$ and complex $`U`$. For $`s=(2+ϵ)`$ this allows us to evaluate the Epstein function up to order $`𝒪(ϵ^2)`$, giving an expression for $`_{(2+ϵ)}`$ up to $`𝒪(ϵ)`$. We use that
$`E_1[z;s;\tau ,c_1]`$ $``$ $`{\displaystyle \underset{n_20}{}}[z+\tau (n_2+c_1)^2]^s`$ (C-2)
has the asymptotic expansion
$`E_1[z;s;\tau ,c_1]`$ $``$ $`z^s{\displaystyle \underset{m0}{}}{\displaystyle \frac{\mathrm{\Gamma }[s+m]}{m!\mathrm{\Gamma }[s]}}\left[{\displaystyle \frac{\tau }{z}}\right]^m\zeta [2m,c_1]+{\displaystyle \frac{z^{1/2s}}{2}}\left[{\displaystyle \frac{\pi }{\tau }}\right]^{\frac{1}{2}}{\displaystyle \frac{\mathrm{\Gamma }[s1/2]}{\mathrm{\Gamma }[s]}}`$ (C-3)
$`+`$ $`{\displaystyle \frac{2\pi ^s}{\mathrm{\Gamma }[s]}}\tau ^{\frac{s}{2}\frac{1}{4}}z^{\frac{s}{2}+\frac{1}{4}}{\displaystyle \underset{p1}{}}p^{s\frac{1}{2}}\mathrm{cos}(2\pi c_1p)K_{s\frac{1}{2}}\left(2\pi p\left(z/\tau \right)^{\frac{1}{2}}\right)`$
Here $`\zeta [q,a],a0,1,2,\mathrm{}`$ is the Hurwitz Zeta function, $`K_s`$ is the Bessel function (E-2). In the $`m=0`$ term in (C-3) one can use $`\zeta [0,c_1]=1/2c_1`$ even for $`c_1=0`$. One can use (C-3) recurrently for the 2D case . With the substitutions
$`c_1c_1(n_1)\rho _2U_1(n_1+\rho _1);zz(n_1)\tau U_2^2(n_1+\rho _1)^2+\delta `$ (C-4)
in eq.(C-3) and after applying a summation over $`n_1`$, one obtains from (C-3) $`_s`$ of (C-1)
$`_s`$ $`=`$ $`+𝒞+𝒟`$
$``$ $`=`$ $`\pi ^s\left[{\displaystyle \frac{\pi }{\tau }}\right]^{\frac{1}{2}}{\displaystyle \frac{\mathrm{\Gamma }[s1/2]}{2}}{\displaystyle \underset{n_1<0}{}}z(n_1)^{\frac{1}{2}s}`$
$`𝒞`$ $`=`$ $`2\tau ^{\frac{s}{2}\frac{1}{4}}{\displaystyle \underset{n_1<0;p1}{}}z(n_1)^{\frac{s}{2}+\frac{1}{4}}p^{s\frac{1}{2}}\mathrm{cos}\left[2\pi c_1(n_1)p\right]K_{s\frac{1}{2}}\left(2\pi p\left(z(n_1)/\tau \right)^{\frac{1}{2}}\right)`$
$`𝒟`$ $`=`$ $`\pi ^s{\displaystyle \underset{n_1<0,m0}{}}{\displaystyle \frac{\mathrm{\Gamma }[s+m]}{m!}}(\tau )^m\zeta [2m,c_1(n_1)]z(n_1)^{sm}`$ (C-5)
The series in $`𝒟`$ is asymptotic . To compute $``$ one considers the cases $`\delta /(\tau U_2^2)<1`$, $`>1`$. In the following we take $`s=2ϵ`$.
If $`0\delta /(\tau U_2^2)<1`$, one uses for $``$ a binomial expansion of its term $`z(n_1)^{5/2+ϵ}`$, as in eqs.(B-7), (B), and the comments thereafter to isolate the poles, to find
$``$ $`=`$ $`{\displaystyle \frac{1}{2}}\{{\displaystyle \frac{\pi ^3\delta ^3}{12\tau |U_2|}}[{\displaystyle \frac{1}{ϵ}}+\mathrm{ln}\left(4\pi \tau U_2^2e^{\gamma +2\psi (1\rho _1)}\right)]+\pi ^3\delta ^2|U_2|[{\displaystyle \frac{1}{12}}{\displaystyle \frac{1}{2}}\rho _1(1\rho _1)]`$ (C-6)
$`+{\displaystyle \frac{2\pi ^3}{45}}\tau ^2|U_2|^5\left[{\displaystyle \frac{1}{21}}\rho _1^2(15\rho _1^22\rho _1^4+6\rho _1^3)\right]{\displaystyle \frac{\pi ^3\delta \tau |U_2|^3}{90}}\left[130\rho _1^2(1\rho _1)^2\right]`$
$`+\pi ^{5/2}\tau ^2|U_2|^5{\displaystyle \underset{p1}{}}{\displaystyle \frac{\mathrm{\Gamma }[p+1/2]}{(p+3)!}}\left[{\displaystyle \frac{\delta }{\tau U_2^2}}\right]^{p+3}\zeta [2p+1,1\rho _1]\}+𝒪(ϵ)`$
If instead $`\delta /(\tau U_2^2)1`$, one uses for $`z(n_1)^{5/2+ϵ}`$ in $``$ the expansion eq.(C-3), to find
$``$ $`=`$ $`{\displaystyle \frac{1}{2}}\{{\displaystyle \frac{\pi ^3\delta ^3}{12\tau |U_2|}}[{\displaystyle \frac{1}{ϵ}}+\mathrm{ln}\left(\pi \delta e^{\gamma \frac{11}{6}}\right)]+2\sqrt{\tau }U_2^2\delta ^{\frac{3}{2}}{\displaystyle \underset{\stackrel{~}{n}_11}{}}{\displaystyle \frac{\mathrm{cos}[2\pi \stackrel{~}{n}_1\rho _1]}{\stackrel{~}{n}_1^3}}K_3(2\pi \stackrel{~}{n}_1\sqrt{{\displaystyle \frac{\delta }{\tau U_2^2}}})`$ (C-7)
$`+{\displaystyle \frac{(\pi \delta )^{5/2}}{\sqrt{\tau }}}{\displaystyle \underset{m0}{}}{\displaystyle \frac{\mathrm{\Gamma }[m5/2]}{m!}}\left({\displaystyle \frac{\tau U_2^2}{\delta }}\right)^m\zeta [2m,1\rho _1]\}+𝒪(ϵ)`$
For $`𝒞`$ one uses the definition of Bessel functions $`K_{5/2}`$ and of $`\text{Li}_\sigma `$, to find
$`𝒞={\displaystyle \frac{\tau }{2}}{\displaystyle \underset{n_1<0}{}}[z(n_1)\text{Li}_3(e^{2\pi \gamma (n_1)})+{\displaystyle \frac{3}{2\pi }}(\tau z(n_1))^{\frac{1}{2}}\text{Li}_4(e^{2\pi \gamma (n_1)})+{\displaystyle \frac{3\tau }{4\pi ^2}}\text{Li}_5(e^{2\pi \gamma (n_1)})+c.c.]+𝒪(ϵ)`$
$`\gamma (n_1)`$ $``$ $`{\displaystyle \frac{1}{\sqrt{\tau }}}(\delta +\tau U_2^2(n_1+\rho _1)^2)^{1/2}i(\rho _2U_1(n_1+\rho _1))`$ (C-8)
Therefore if $`0\delta /(\tau U_2^2)1`$ one has from (C), (C-6), (C)
$`_{(2+ϵ)}`$ $`=`$ $`𝒟+{\displaystyle \frac{(\tau U_2)^2}{2}}\{{\displaystyle \frac{\pi ^3\delta ^3}{12(\tau |U_2|)^3}}{\displaystyle \frac{1}{ϵ}}+{\displaystyle \frac{2\pi ^3}{45}}|U_2|^3[{\displaystyle \frac{1}{21}}\rho _1^2(15\rho _1^22\rho _1^4+6\rho _1^3)]`$ (C-9)
$`+`$ $`{\displaystyle \underset{n_1<0}{}}[(n_1+\rho _1)^2\text{Li}_3(\sigma _{n_1})+{\displaystyle \frac{3|n_1+\rho _1|}{2\pi |U_2|}}\text{Li}_4(\sigma _{n_1})+{\displaystyle \frac{3}{4\pi ^2U_2^2}}\text{Li}_5(\sigma _{n_1})+c.c.]\}`$
$`\sigma _{n_1<0}=e^{2i\pi \left(\rho _2U_1\left|\rho _1+n_1\right|\right)}e^{2\pi |U_2|(\rho _1+n_1)},0\rho _{1,2}<1`$
Eq.(3.2) in the text is then $`V_𝒪=\mu ^4/(2\pi )^{4ϵ}_{(2\frac{ϵ}{2})}\left(\delta m^2/\mu ^2,\tau (2\pi )^2/(\mu ^2𝒜U_2)\right)`$.
It remains to evaluate the series of $`𝒟`$ in (C) in a form amenable to numerical evaluation. For this, its factor $`z(n_1)`$ under the sum can be expanded in $`\delta /(\tau U_2^2)`$ if $`0\delta /(\tau U_2^2)<1`$ by using the binomial expansion eq.(B-7). If $`\delta /(\tau U_2^2)>1`$ one uses instead the asymptotic expansion of eq.(C-3). For our purposes $`s=(2+ϵ)`$, and then
$$𝒟=\pi ^{2+ϵ}\underset{n_1<0,m0}{}\mathrm{\Gamma }[m2ϵ]\frac{(\tau )^m}{m!}\zeta [2m,\rho _2U_1(n_1+\rho _1)]\left[\delta +\tau U_2^2(n_1+\rho _1)^2\right]^{2m+ϵ}$$
(C-10)
Eq.(48) in the text is $`\stackrel{~}{𝒟}2𝒜^2\mu ^4/(2\pi )^{4ϵ}𝒟\left(ϵϵ/2,\tau (2\pi )^2/(\mu ^2U_2𝒜),\delta m^2/\mu ^2\right)`$.
$``$ Case of $`𝒯_2/Z_4`$ orbifold: In the following we restrict the calculation of $`𝒟`$ to $`𝒯_2/Z_4`$, when $`U_1=0`$, $`U_2=1`$. If so, the argument of zeta function in $`𝒟`$ does not have a $`n_1`$ dependence and the sums over $`n_1`$ and $`m`$ can be easily performed. (For other orbifolds $`U_10`$ further evaluation of $`𝒟`$ is more tedious but very similar).
(a). For $`0\delta /\tau <1`$, after a binomial expansion (B-7) of last bracket in (C-10), $`𝒟`$ becomes
$`𝒟`$ $`=`$ $`(\pi \tau )^{2+ϵ}{\displaystyle \underset{n_1<0,m0}{}}\mathrm{\Gamma }[m2ϵ]{\displaystyle \frac{(1)^m}{m!}}\zeta [2m,\rho _2]\left[\delta /\tau +(n_1+\rho _1)^2\right]^{2m+ϵ}`$ (C-11)
$`=`$ $`(\pi \tau )^{2+ϵ}{\displaystyle \underset{m0,p0}{}}\left[{\displaystyle \frac{\delta }{\tau }}\right]^p(1)^{m+p}{\displaystyle \frac{\mathrm{\Gamma }[p+m2ϵ]}{m!p!}}\zeta [2m,\rho _2]\zeta [2p+2m42ϵ,1\rho _1]`$
$`=`$ $`{\displaystyle \frac{(\pi \tau )^2}{2ϵ}}\left\{\zeta [4,\rho _2]\zeta [0,\rho _1]+\zeta [4,\rho _1]\zeta [0,\rho _2]+2\zeta [2,\rho _1]\zeta [2,\rho _2]\right\}`$
$`+`$ $`{\displaystyle \frac{\pi ^2\delta \tau }{ϵ}}\left\{\zeta [0,\rho _1]\zeta [2,\rho _2]+\zeta [2,\rho _1]\zeta [0,\rho _2]\right\}+{\displaystyle \frac{\pi ^2\delta ^2}{2ϵ}}\zeta [0,\rho _1]\zeta [0,\rho _2]+𝒟_f+𝒪(ϵ)`$
used in (3.2.1), (53) with $`\stackrel{~}{𝒟}2\mu ^4L^4/(2\pi )^{4ϵ}𝒟\left(ϵϵ/2,\tau (2\pi )^2/(\mu L)^2,\delta m^2/\mu ^2\right)`$.
$`𝒟_f`$ in eq.(C-11) is the finite $`𝒪(ϵ^0)`$ part:
$`𝒟_f`$ $`=`$ $`{\displaystyle \frac{(\pi \tau )^2}{2}}[2\zeta [2,\rho _1]\zeta [2,\rho _2]\mathrm{ln}(\tau \pi e^{\gamma 1})+\zeta [0,\rho _2](\zeta [4,\rho _1]\mathrm{ln}(\pi \tau e^{\gamma \frac{3}{2}})+2\zeta ^{}[4,1\rho _1])`$ (C-12)
$`+`$ $`4\zeta [2,\rho _2]\zeta ^{}[2,1\rho _1]+\zeta [4,\rho _2](\zeta [0,\rho _1]\mathrm{ln}(\pi \tau e^\gamma )+2\zeta ^{}[0,1\rho _1])]`$
$`+`$ $`\pi ^2\delta \tau [\zeta [0,\rho _2](\zeta [2,\rho _1]\mathrm{ln}(\pi \tau e^{\gamma 1})+2\zeta ^{}[2,1\rho _1])+\zeta [2,\rho _2](\zeta [0,\rho _1]\mathrm{ln}(\pi \tau e^\gamma )`$
$`+`$ $`2\zeta ^{}[0,1\rho _1])]+{\displaystyle \frac{(\pi \delta )^2}{2}}\zeta [0,\rho _2](\zeta [0,\rho _1]\mathrm{ln}(\pi \tau e^\gamma )+2\zeta ^{}[0,1\rho _1])`$
$`+`$ $`{\displaystyle \underset{p0,m0,p+m3}{}}\pi ^2\tau ^2\left({\displaystyle \frac{\delta }{\tau }}\right)^p(1)^{p+m}{\displaystyle \frac{\mathrm{\Gamma }[p+m2]}{p!m!}}\zeta [2m,\rho _2]\zeta [2p+2m4,1\rho _1]`$
This was used in (3.2.2), (57) with $`\stackrel{~}{𝒟}_f2\mu ^4L^4/(2\pi )^4𝒟_f\left(\tau (2\pi )^2/(\mu L)^2,\delta m^2/\mu ^2\right)`$, after neglecting any $`mL1`$ dependence.
(b). In the case when $`\delta /\tau >1`$ one uses in $`𝒟`$ of (C-10) or the first line in (C-11), the asymptotic expansion eq.(C-3). The results shows that the divergent part of $`𝒟`$ is identical to that in the last two lines in (C-11), while the value of $`𝒟_f`$ ($`𝒪(ϵ^0)`$) in (C-11) has now the expression
$`𝒟_f`$ $`=`$ $`(\pi \tau )^2{\displaystyle \underset{m0,k0,k+m3}{}}(1)^{m+k}\zeta [2k,1\rho _1]\zeta [2m,\rho _2]{\displaystyle \frac{\mathrm{\Gamma }[k+m2]}{m!k!}}\left({\displaystyle \frac{\delta }{\tau }}\right)^{2mk}`$
$`+`$ $`{\displaystyle \underset{m0}{}}{\displaystyle \frac{\pi ^{\frac{5}{2}}\tau ^2}{2}}\zeta [2m,\rho _2]{\displaystyle \frac{(1)^m}{m!}}\left({\displaystyle \frac{\delta }{\tau }}\right)^{\frac{5}{2}m}\mathrm{\Gamma }[m5/2]`$
$`+`$ $`{\displaystyle \underset{m0}{}}{\displaystyle \frac{2\tau ^2(\pi )^m}{m!}}\left({\displaystyle \frac{\delta }{\tau }}\right)^{\frac{5}{4}\frac{m}{2}}\zeta [2m,\rho _2]{\displaystyle \underset{p1}{}}\mathrm{cos}[2\pi p\rho _1]p^{m\frac{5}{2}}K_{m\frac{5}{2}}\left(2\pi p\sqrt{{\displaystyle \frac{\delta }{\tau }}}\right)`$
$`+`$ $`{\displaystyle \frac{\pi ^2\tau ^2}{2}}\left[\zeta [4,\rho _2]\zeta [0,\rho _1]+\zeta [4,\rho _1]\zeta [0,\rho _2]+2\zeta [2,\rho _1]\zeta [2,\rho _2]\right]\mathrm{ln}(\pi \delta e^\gamma )`$
$`+`$ $`\pi ^2\tau \delta (\zeta [2,\rho _1]\zeta [0,\rho _2]+\zeta [2,\rho _2]\zeta [0,\rho _1])\mathrm{ln}(\pi \delta e^{\gamma 1})+{\displaystyle \frac{\pi ^2\delta ^2}{2}}\zeta [0,\rho _1]\zeta [0,\rho _2]\mathrm{ln}(\pi \delta e^{\gamma \frac{3}{2}})`$
To conclude, if $`\delta /\tau 1`$, the value of $`_{(2+ϵ)}`$ is given by
$`_{(2+ϵ)}`$ $`=`$ $`+𝒞+𝒟_f`$ (C-14)
$`+`$ $`{\displaystyle \frac{(\pi \tau )^2}{2ϵ}}\left\{\zeta [4,\rho _2]\zeta [0,\rho _1]+\zeta [4,\rho _1]\zeta [0,\rho _2]+2\zeta [2,\rho _1]\zeta [2,\rho _2]\right\}`$
$`+`$ $`{\displaystyle \frac{\pi ^2\delta \tau }{ϵ}}\left\{\zeta [0,\rho _1]\zeta [2,\rho _2]+\zeta [2,\rho _1]\zeta [0,\rho _2]\right\}+{\displaystyle \frac{\pi ^2\delta ^2}{2ϵ}}\zeta [0,\rho _1]\zeta [0,\rho _2]+𝒪(ϵ)`$
$``$ is given in eq.(C-7), $`𝒞`$ in eq.(C), while $`𝒟_f`$ is that of eq.(C).
Eq.(C-14) concludes the calculation of $`_{(2+ϵ)}`$ for $`𝒯_2/Z_4`$ for $`\delta /\tau 1`$. Eq.(C-9) with (C-11), (C-12), gives $`_{(2+ϵ)}`$ for $`0<\delta /\tau <1`$ again for $`𝒯_2/Z_4`$. Eq.(C-9) with (C-6), (C-7), (C), (C-10) give $`_{(2+ϵ)}`$ for any $`U`$.
With the expression (C-14) for $``$, eq.(60) in the text is then
$`V_𝒪={\displaystyle \frac{\mu ^4}{(2\pi )^{4ϵ}}}_{(2\frac{ϵ}{2})}\left(\delta {\displaystyle \frac{m^2}{\mu ^2}},\tau {\displaystyle \frac{(2\pi )^2}{(\mu L)^2}}\right).`$ (C-15)
### D . Orbifolds, Fixed points and discrete Wilson lines.
The lattice of $`T_2/Z_N`$ orbifolds is generated by $`(1,\xi )`$ with $`\xi =i`$ for $`Z_2`$, $`Z_4`$ and $`\xi =e^{2i\pi /3}`$ for $`Z_3`$, $`Z_6`$. The group $`Z_N`$ of discrete rotations has $`N`$ elements $`P_N^n`$, $`0n<N1`$ with $`P_N^N=1`$. Their fixed points are
$`Z_2:`$ $`P_2:z_{f.p.}=0,\mathrm{\hspace{0.17em}\hspace{0.17em}1}/2,\xi /2,(1+\xi )/2,\xi =i`$
$`Z_3:`$ $`P_3,P_3^2:z_{f.p.}=0,(2+\xi )/3,(1+2\xi )/3,\xi =e^{2i\pi /3}`$
$`Z_4:`$ $`P_4,P_4^3:z_{f.p.}=0,(1+\xi )/2,\xi =i`$
$`P_4^2:z_{f.p.}=0,(1+\xi )/2,\xi /2,\mathrm{\hspace{0.17em}\hspace{0.17em}1}/2.`$
$`Z_6:`$ $`P_6,P_6^5:z_{f.p.}=0.\xi =e^{2i\pi /3}`$ (D-1)
$`P_6^2,P_6^4:z_{f.p.}=0,(2+\xi )/3,(1+2\xi )/3,`$
$`P_6^3:z_{f.p.}=0,\mathrm{\hspace{0.17em}\hspace{0.17em}1}/2,\xi /2,(1+\xi )/2.`$
The usual orbifold action ($`g`$) and that of Wilson lines ($`T_{1,2}`$) are given by
$`\mathrm{\Phi }^g(z+1)`$ $`=`$ $`T_1\mathrm{\Phi }^g(z),`$
$`\mathrm{\Phi }^g(z+\tau )`$ $`=`$ $`T_2\mathrm{\Phi }^g(z)`$
$`\mathrm{\Phi }^g(\tau z)`$ $`=`$ $`g\mathrm{\Phi }^g(z),\mathrm{with}\tau e^{2i\pi /N}`$ (D-2)
One has that
$`\mathrm{\Phi }^g(z+\tau )=\mathrm{\Phi }^g(\tau (\tau ^1z+1))=gT_1\mathrm{\Phi }^g(\tau ^1z)=gT_1g^{N1}\mathrm{\Phi }^g(z)`$ (D-3)
Using the definition of $`T_2`$, then
$`T_2g=gT_1`$ (D-4)
One can further assume that the orbifold action $`g`$ and the Wilson lines $`T_i`$ commute, then $`T_1=T_2=T`$, and with $`T_ie^{2i\pi \rho _i}`$ one finds (modulo $`𝐙`$) that $`\rho _1=\rho _2=\rho `$ for any $`T_2/Z_N`$.
The case of $`Z_3`$ orbifolds: ($`\xi =e^{2i\pi /3}`$)
$`\mathrm{\Phi }^g(\tau ^2(z+1))`$ $`=`$ $`g^2T_1\mathrm{\Phi }^g(z)`$
$`\mathrm{\Phi }^g(\tau ^2(z+1))`$ $`=`$ $`T_1^1\mathrm{\Phi }^g(\tau ^2(z+1)+1)=T_1^1T_2^1\mathrm{\Phi }^g(\tau ^2z)=T_1^1T_2^1g^2\mathrm{\Phi }^g(z)`$ (D-5)
A solution to this is
$`g^2T_1=T_1^1T_2^1g^2`$ (D-6)
or, assuming $`[g,T_i]=0`$ giving $`T_1=T_2=T`$
$`T^3=1,\rho =0,\mathrm{\hspace{0.17em}1}/3,\mathrm{\hspace{0.17em}2}/3.`$ (D-7)
where $`T=\mathrm{exp}(2i\pi \rho )`$. Further, for the fixed points
$`a).ifz_{f.p.}`$ $`=`$ $`0,\mathrm{\Phi }^g(0)=\mathrm{\Phi }^g(\tau \mathrm{\hspace{0.17em}0})=g\mathrm{\Phi }^g(0)`$
$`b).ifz_{f.p.}`$ $`=`$ $`(2+\xi )/3,Tg\mathrm{\Phi }^g(z_{f.p.})=\mathrm{\Phi }^g(z_{f.p.})`$
$`c).ifz_{f.p.}`$ $`=`$ $`(1+2\xi )/3,T^2g\mathrm{\Phi }^g(z_{f.p.})=\mathrm{\Phi }^g(z_{f.p.})=Tg^2\mathrm{\Phi }^g(z_{f.p.})`$ (D-8)
These are additional conditions which must be respected by the Wilson lines $`T`$, orbifold projections $`g`$ and fields $`\mathrm{\Phi }^g`$ at the fixed points. The conditions can be respected by suitable relative choices for $`T`$, $`g`$, or trivially by requiring the fields vanish at these fixed points.
The case of $`Z_4`$ orbifolds: ($`\xi =i`$)
$`\mathrm{\Phi }^g(\tau ^2(z+1))`$ $`=`$ $`T_1^1\mathrm{\Phi }^g(\tau ^2z+\tau ^2+1)=T_1^1g^2\mathrm{\Phi }^g(z)`$
$`\mathrm{\Phi }^g(\tau ^2(z+1))`$ $`=`$ $`g^2T_1\mathrm{\Phi }^g(z)`$ (D-9)
which gives
$`g^2T_1=T_1^1g^2`$ (D-10)
or, assuming $`[g,T_i]=0`$ giving $`T_1=T_2=T`$ one has
$`T^2=1,\rho =0,1/2.`$ (D-11)
where $`T=\mathrm{exp}(2i\pi \rho )`$. Further, for the fixed points
$`a).ifz_{f.p.}`$ $`=`$ $`0,\mathrm{\Phi }^g(0)=\mathrm{\Phi }^g(\tau \mathrm{\hspace{0.17em}0})=g\mathrm{\Phi }^g(0)`$
$`b).ifz_{f.p.}`$ $`=`$ $`(1+\xi )/2,Tg\mathrm{\Phi }^g(z_{f.p.})=\mathrm{\Phi }^g(z_{f.p.})`$
$`c).ifz_{f.p.}`$ $`=`$ $`1/2,Tg^2\mathrm{\Phi }^g(z_{f.p.})=T\mathrm{\Phi }^g(1/2)=\mathrm{\Phi }^g(z_{f.p.})`$
$`d).ifz_{f.p.}`$ $`=`$ $`\xi /2,Tg^2\mathrm{\Phi }^g(z_{f.p.})=gT\mathrm{\Phi }^g(\xi ^2/2)=\mathrm{\Phi }^g(z_{f.p.})`$ (D-12)
Similar to the $`Z_3`$ case, the conditions can be respected by suitable choices for $`T`$, $`g`$, or trivially by requiring the fields vanish at these fixed points.
The case of $`Z_6`$ orbifolds:
$`\mathrm{\Phi }^g(\tau ^3(z+1))`$ $`=`$ $`g^3T_1\mathrm{\Phi }^g(z)`$
$`\mathrm{\Phi }^g(\tau ^3(z+1))`$ $`=`$ $`T_1^1\mathrm{\Phi }^g(\tau ^3z+\tau ^3+1)=T_1^1g^3\mathrm{\Phi }^g(z)`$ (D-13)
since $`\tau ^3+1=0`$. One solution is
$`g^3T_1=T_1^1g^3`$ (D-14)
Assuming $`[g,T_1]=0`$ which gives $`T_1=T_2=T`$ one has
$`T^2=1,\rho =0,1/2.`$ (D-15)
where $`T\mathrm{exp}(2i\pi \rho )`$. Further,
$`\mathrm{\Phi }^g(\tau ^2(z+1))`$ $`=`$ $`g^2T_1\mathrm{\Phi }^g(z)`$
$`\mathrm{\Phi }^g(\tau ^2(z+1))`$ $`=`$ $`T_1^1\mathrm{\Phi }^g(\tau ^2z+\tau ^2+1)=T_1^1T_2g^2\mathrm{\Phi }^g(z)`$ (D-16)
since $`\tau ^2\tau +1=0`$. One solution is
$`g^2T_1=T_1^1T_2g^2`$ (D-17)
With $`[g,T_i]=0`$ giving $`T_1=T_2=T`$ one has
$`T=1\rho =0.`$ (D-18)
where $`T\mathrm{exp}(2i\pi \rho )`$. Thus, if $`[g,T_i]=0`$, one concludes from (D-15), (D-18) that $`\rho =0`$. Further relations at the fixed points exist, which can be found as in the case of $`𝒯_2/Z_{3,4}`$.
### E . Mathematical Formulae and Conventions.
We used the Poisson re-summation formula
$$\underset{nZ}{}e^{\pi A(n+\sigma )^2}=\frac{1}{\sqrt{A}}\underset{\stackrel{~}{n}Z}{}e^{\pi A^1\stackrel{~}{n}^2+2i\pi \stackrel{~}{n}\sigma }$$
(E-1)
The integral representation of Bessel Function $`K_\nu `$
$$_0^{\mathrm{}}𝑑xx^{\nu 1}e^{bx^pax^p}=\frac{2}{p}\left[\frac{a}{b}\right]^{\frac{\nu }{2p}}K_{\frac{\nu }{p}}(2\sqrt{ab}),Re(b),Re(a)>0$$
(E-2)
with
$`K_{\frac{5}{2}}(z)`$ $`=`$ $`\sqrt{{\displaystyle \frac{\pi }{2z}}}e^z\left[1+{\displaystyle \frac{3}{z^2}}+{\displaystyle \frac{3}{z^2}}\right]`$
$`K_3(z1)`$ $`=`$ $`e^z\sqrt{{\displaystyle \frac{\pi }{2z}}}\left[1+{\displaystyle \frac{35}{8}}{\displaystyle \frac{1}{z}}+{\displaystyle \frac{945}{128}}{\displaystyle \frac{1}{z^2}}+{\displaystyle \frac{3465}{1024}}{\displaystyle \frac{1}{z^3}}+\mathrm{}\right]`$ (E-3)
The definition of PolyLogarithm
$`Li_\sigma (x)={\displaystyle \underset{n1}{}}{\displaystyle \frac{x^n}{n^\sigma }}`$ (E-4)
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# Magnon transport and spin current switching through quantum dots
## I Introduction
Electronic transport in mesoscopics has been widely studied for the last two decades. Particularly due to the interplay between the electronic charge and spin many novel features have been found, which lead to new inventions as giant magnetoresistance, tunneling magnetoresistance or the Datta-Das transistor (see e.g. Rfs. Prinz, 1995; Awschalom and Kikkawa, 1999; Wolf et al., 2001; S. Maekawa, T. Shinjo, 2002 for reviews). This has opened the field of spintronics with possible future applications in the field of quantum computing. The fact that spin relaxation and dephasing rates are rather long compared to charge excitations life-times has recently stimulated the investigation of spin transport in purely magnetic systems without an accompanying charge transport. Experimentally, anomalous heat transport has been reported in $`(Sr,Ca,La)_{14}Cu_{24}O_{41}`$-materials Sologubenko et al. (2000), which is explained by the specific dynamics of spin excitations in quasi-onedimensional systems Zotos et al. (1997). These experiments are still in the incoherent regime where the mean free path of the spinons is much smaller than the system size. In contrast, Meier and LossF. Meier and D. Loss (2003) investigated coherent spin transport through mesoscopic ferromagnetic and antiferromagnetic Heisenberg spin chains. For the case of the isotropic antiferromagnetic spin-$`1/2`$ chain they found that the spin conductance is quantized in units of order $`(g\mu _\mathrm{B})^2/h`$. MatveevMatveev (2004) studied spinon transport through a one-dimensional Wigner crystal, which effectively can be described by an antiferromagnetically ordered spin chain. Here, anomalous conductance quantization phenomena occur which might explain experiments through quantum point contacts Thomas et al. (1996).
Motivated by these interesting recent results for spin transport through one-dimensional spin chains, we analyze in this work the magnetic analog of electronic transport through zero-dimensional quantum dots in the weak coupling regime (i.e. weak exchange coupling between the magnetic quantum dot and the magnetization reservoirs). In particular, we study the spin current through a quantum dot consisting of two interacting spin-$`1/2`$, coupled to two ferromagnetic insulators, which are for simplicity described by a ferromagnetic Heisenberg model (the basic transport features revealed here are expected to hold for any reservoir involving magnon-like excitations). For the same system Wang et al.B. Wang, J. Wang, J. Wang and D.Y. Xing (2004) derived a Landauer-Büttiker-type formula for spin current transport. However, they studied the magnon transport using a magnon representation for the dot states and treated the intra-dot magnon-magnon interaction in a mean-field approximation. In contrast, we study the case of arbitrarily strong intra-dot interaction and, using a real-time formalism developed in Ref. H. Schoeller, 1997, perform only a perturbative expansion in the exchange coupling between the quantum dot and the magnetization reservoirs. This is the regime where in electronic transport interesting effects like Coulomb blockade phenomena occur.
Our main focus lies on the investigation of possibilities to switch the spin current by tuning system parameters as e.g. an intra-dot magnetic field or the exchange coupling between the two spin-$`1/2`$. In Rfs. G. Burkard, D. Loss, D.P. DiVincenzo, 1999; G. Burkard, G. Seelig, D. Loss, 2000 it was shown that the latter may be modulated, e.g. by application of gates. In contrast to electronic transport, we find that the specific microscopic coupling configuration and the magnon dispersion relation in the reservoirs has a strong influence on the spin current. Interestingly, the spin current can not only be dominated by 1-magnon transport processes, but also by 2-magnon cycles, leading to specific resonances or antiresonances in the spin current as function of the intra-dot exchange coupling. Furthermore, we discuss two different mechanisms of driving the spin current, firstly a difference of left and right magnon chemical potentials, and secondly a magnetic field gradient. We find that they are not equivalent but yield different results, unlike the counterpart of electrochemical potentials in electronic transport.
The paper is organized as follows. In section II we introduce the model and the extension of the real-time formalism to bosonic transport. Assuming the reservoirs to be magnetized in the same direction, subsequently, in section III.1 the spin current is discussed for the case of differing chemical potentials, in section III.2 for the case of a magnetic field gradient applied, in each case for two different dot-reservoir coupling configurations. In section III.3 we deal with the case of antiparallel magnetized reservoirs and finally close with some continuative remarks concerning extensions to this work.
## II Model
Our model consists of a quantum dot, coupled via local exchange to two magnetization reservoirs. The dot is made up of two coupled spin-$`1/2`$, whereas the reservoirs are assumed to be ferromagnetic insulators, which can be described by a ferromagnetic Heisenberg model, see Fig. 2.
The total Hamiltonian $`\overline{}`$ for the system can be written as sum of the dot Hamiltonian, the reservoir contributions and the exchange part, describing the coupling between the dot and the reservoirs (in the following called tunneling part corresponding to tunneling of magnons instead of electrons in the electronic analog):
$$\overline{}=\overline{}_{\mathrm{d}\mathrm{o}\mathrm{t}}+\overline{}_L+\overline{}_R+\overline{}_T.$$
(1)
The dot Hamiltonian $`\overline{}_{\mathrm{d}\mathrm{o}\mathrm{t}}`$ reads
$$\overline{}_{\mathrm{d}\mathrm{o}\mathrm{t}}=J_{\mathrm{d}\mathrm{o}\mathrm{t}}𝐬_1𝐬_2+B_{\text{Zee}}s^z,$$
(2)
where $`J_{\mathrm{d}\mathrm{o}\mathrm{t}}`$ is the exchange coupling between the two spin-$`1/2`$, $`s^z=s_1^z+s_2^z`$ and $`B_{\text{Zee}}`$ represents a homogeneous magnetic field applied on the dot.<sup>1</sup><sup>1</sup>1 The factor $`g\mu _\mathrm{B}`$ is included in the Zeeman energy $`B_{\text{Zee}}`$, in the following referred to as a magnetic field. This is also the case with all following magnetic fields, resp. Zeeman energies. The reservoir contributions are given by a Heisenberg Hamiltonian:
$$\overline{}_r=J/2\underset{i,j}{}𝐒_{ri}𝐒_{rj}+B_r\underset{i}{}S_{ri}^z+b_1S_{r1}^z,$$
(3)
with $`r\{L,R\}`$, $`J>0`$ and $`i,j`$ denoting the sum over neighboring sites.<sup>2</sup><sup>2</sup>2 In this work we assume one-dimensional reservoirs. However, since the qualitative shape of the spectral function does not depend on the reservoir dimension, apart from minor modulations, no qualitative changes occur for reservoirs of other dimensionality. $`B_r`$ represents a homogeneous magnetic field applied on the reservoir $`r`$, whereas $`b_1`$ models a remnant of the dot magnetic field $`B_{\text{Zee}}`$. For simplicity we assume that this remnant field acts exclusively on the reservoir sites neighboring the dot.<sup>3</sup><sup>3</sup>3 The consideration of the remnant field $`b_1`$ is necessary to break isotropy. This is due to a subtlety occurring in the calculation of the spectral function $`\mathrm{\Gamma }(\omega )`$ \[cf. (18)\] for the semi-infinite Heisenberg model. As a consequence of this $`b_1=0`$ represents an unphysical border case.
Isotropy may also be broken by considering an anisotropic Heisenberg model (z-anisotropy). However, as long as a ferromagnetic ground state prevails, this does not lead to qualitative modifications of the results presented here. The tunneling part is given by
$$\overline{}_T=\underset{r=L,R}{}\underset{i=1,2}{}J_r^i𝐬_i𝐒_{r1}.$$
(4)
where $`J_r^i`$ are the real exchange couplings between the two dot spins and the adjacent spins of the reservoirs, in the following referred to as tunnel couplings.
For the case of parallel magnetized reservoirs we assume the reservoir ground states to be $`|S_i^z=S,i`$ (the modifications for the case of antiparallel magnetized reservoirs will be stated in section III.3). Making use of the Holstein-Primakoff transformation, $`\overline{}_r`$ may be diagonalized by introducing the magnon creation (annihilation) operators $`a_{rk}^{}`$ ($`a_{rk}`$):
$$_r=\underset{k}{}\omega _{rk}a_{rk}^{}a_{rk},\omega _{rk}=2JS[1\mathrm{cos}(ka)]+B_r,$$
(5)
with the lattice constant $`a`$. Here and throughout the paper we set $`\mathrm{}=k_\mathrm{B}=1`$. Furthermore, in the following all energies as also the spin current are normalized to $`JS`$.
Rewriting $`\overline{}_T`$ in terms of the reservoir magnon creation and annihilation operators yields
$$_T=\underset{r,i,k}{}\left(J_{rk}^is_i^+a_{rk}+\text{ }\text{h.\hspace{0.17em}c.}\right),$$
(6)
with the $`k`$-dependent coupling $`J_{rk}^iJ_r^i1|k`$, where $`1|k`$ represents the amplitude of the magnon wavefunction at the dot neighboring reservoir sites. In (6) we have split off the $`s_i^zS_{r1}^z`$-term arising in (4) and have included its ground state average (i.e. replacing $`S_{r1}^zS`$) into the dot Hamiltonian. This leads to an effective field $`H_i=_rJ_r^iS`$ for the dot states:
$$_{\mathrm{d}\mathrm{o}\mathrm{t}}=J_{\mathrm{d}\mathrm{o}\mathrm{t}}𝐬_1𝐬_2+B_{\text{Zee}}s^z+\left(H_1s_1^z+H_2s_2^z\right).$$
(7)
Thus our final Hamiltonian $``$ reads
$$=_0+_T,_0=_L+_R+_{\mathrm{d}\mathrm{o}\mathrm{t}}.$$
(8)
The triplet states $`|𝖳_+=|`$ and $`|𝖳_{}=|`$ are eigenstates of $`_{\mathrm{d}\mathrm{o}\mathrm{t}}`$. Within the sub-space spanned by $`|𝖲=\frac{1}{\sqrt{2}}\left(\right||)`$ and $`|𝖳_\mathrm{𝟢}=\frac{1}{\sqrt{2}}\left(\right|+|)`$, i.e. the states with $`s^z`$ quantum number $`m=0`$, $`_{\mathrm{d}\mathrm{o}\mathrm{t}}`$ corresponds to the matrix $`\left[\begin{array}{cc}J_{\mathrm{d}\mathrm{o}\mathrm{t}}& \mathrm{\Delta }H\\ \mathrm{\Delta }H& 0\end{array}\right]`$. $`\mathrm{\Delta }H=(H_1H_2)/2`$ measures the inhomogeneity of the effective field $`H_i`$, which originates from the tunnel couplings.
We consider two different mechanisms driving the spin current through the dot, which we denote as $`\mathrm{\Delta }\mu `$\- and *$`\mathrm{\Delta }B`$-configuration* in the following.
For the *$`\mathrm{\Delta }\mu `$-configuration* the magnetization of the reservoirs, given by the magnon number, is supposed to be conserved. This means that we assume the spin relaxation time in the reservoirs to be longer than the typical spin current measurement time. Therefore a chemical potential $`\mu _r`$ is assigned to the magnons in each of the reservoirs. Moreover we set $`B_r=0`$ here, so that only the difference between the magnon chemical potentials $`\mu _L`$ and $`\mu _R`$ drives the spin current.
For the *$`\mathrm{\Delta }B`$-configuration* we set $`\mu _r=0`$ in both reservoirs. Here a magnetic field gradient $`B_{L/R}(t)=B_0\pm \mathrm{\Delta }B`$, switched on simultaneously with the tunneling, gives rise to a non-vanishing spin current. The offset field $`B_0|\mathrm{\Delta }B|`$ introduced here is required for the stabilization of the ferromagnetic ground state.
We use the real-time formalism described in Ref. H. Schoeller, 1997 to compute the stationary spin current. By integrating out the reservoir degrees of freedom this formalism yields an effective description of the system in terms of the dot degrees of freedom. It is based on the formally exact kinetic equations (tunneling and magnetic field gradient are assumed to be switched on at time $`t=0`$)
$$I_r(t)=\mathrm{Tr}_{dot}\left[_0^t𝑑t^{}\mathrm{\Sigma }_{I_r}(tt^{})p(t^{})\right],$$
(9)
$$\dot{p}(t)=\text{i}\text{L}_0p(t)+_0^t𝑑t^{}\mathrm{\Sigma }(tt^{})p(t^{}),$$
(10)
where $`I_r(t)`$ denotes the time-dependent expectation value of the spin current through the dot. The current operator $`\widehat{I}_r`$ is given by
$$\widehat{I}_r=\text{i}[,N_r]_{}=\text{i}\underset{i,k}{}\left(J_{rk}^ia_{rk}^{}s_i^{}J_{rk}^is_i^+a_{rk}\right)$$
(11)
with $`[,]_{}`$ denoting the commutator and $`N_r=_ka_{rk}^{}a_{rk}`$, so that the spin current $`I_r`$ is positive, when a ‘spin-up’ leaves the reservoir $`r`$. $`p(t)`$ is the time-dependent reduced density matrix, which is the trace over the reservoir degrees of freedom, $`\mathrm{Tr}_{res}`$, of the full density matrix. In contrast, $`\mathrm{Tr}_{dot}`$ denotes the trace over the local (dot) degrees of freedom. The free propagation of the system is determined by the Liouvillian $`\text{L}_0`$. It is defined by $`\text{L}_0=[_0,]_{}`$. The coupling to the reservoirs is described by the integral kernels $`\mathrm{\Sigma }`$ and $`\mathrm{\Sigma }_{I_r}`$. They are given by
$$\begin{array}{cc}\hfill \mathrm{\Sigma }& (tt^{})\hfill \\ & =(\text{i})^2\mathrm{Tr}_{res}\left[\text{L}_Te^{\text{i}\text{L}_0t}Te^{\text{i}_t^{}^t𝑑\tau \text{L}_T(\tau )}e^{\text{i}\text{L}_0t^{}}\text{L}_T\rho _{\mathrm{r}\mathrm{e}\mathrm{s}}^{\text{eq}}\right]_{\text{irred.}},\hfill \end{array}$$
(12)
$$\begin{array}{cc}\hfill \mathrm{\Sigma }_{I_r}(tt^{})=& \text{i}\mathrm{Tr}_{res}[\text{A}_{I_r}e^{\text{i}\text{L}_0t}Te^{\text{i}_t^{}^t𝑑\tau \text{L}_T(\tau )}\hfill \\ & \times e^{\text{i}\text{L}_0t^{}}\text{L}_T\rho _{\mathrm{r}\mathrm{e}\mathrm{s}}^{\text{eq}}]_{\text{irred.}},\hfill \end{array}$$
(13)
where we assumed that the reservoirs are in equilibrium described by the reservoir density matrix $`\rho _{\mathrm{r}\mathrm{e}\mathrm{s}}^{\text{eq}}`$. $`\text{L}_T=[_T,]_{}`$ is the interaction part of the Liouvillian, and the interaction picture is defined by $`\text{L}_T(t)=e^{\text{i}\text{L}_0t}\text{L}_Te^{\text{i}\text{L}_0t}`$. The superoperator $`\text{A}_{I_r}`$ is given by the anticommutator with the current operator $`\widehat{I}_r`$:
$$\text{A}_{I_r}=\frac{1}{2}[\widehat{I}_r,]_+.$$
(14)
$`T`$ denotes the time ordering operator, and the index “irred.” indicates that only irreducible diagrams are taken into account (meaning that each vertical cut through a diagram will cut at least one bosonic reservoir line, see Ref. H. Schoeller, 1997 for details). Introducing the Laplace transforms $`f(z)=_0^{\mathrm{}}𝑑te^{\text{i}zt}f(t)`$ of the time-dependent functions $`f(t)`$ and using the identity $`lim_t\mathrm{}f(t)=\text{i}lim_{z\text{i}0^+}zf(z)`$ leads to the following expression for the stationary spin current $`I_r^{\text{st}}`$:
$$I_r^{\text{st}}=\mathrm{Tr}_{dot}\left[\mathrm{\Sigma }_{I_r}(\text{i}0^+)p^{\text{st}}\right],$$
(15)
where the stationary reduced density matrix $`p^{\text{st}}`$ is determined by
$$\left[\text{i}\text{L}_0+\mathrm{\Sigma }(\text{i}0^+)\right]p^{\text{st}}=0.$$
(16)
We calculate the kernels $`\mathrm{\Sigma }`$ and $`\mathrm{\Sigma }_{I_r}`$ to lowest order in tunneling, i.e. second order in the coupling constants $`J_r^i`$. The diagrams contributing are listed in the appendix. A typical diagram is shown in Fig. 1(a), which corresponds to the following expression:
$$\begin{array}{cc}\hfill \frac{\text{i}}{2\pi }\underset{i,j}{}s_1|J_{L/R}^i& s_i^+|s_2s_2^{}|J_{L/R}^js_j^{}|s_1^{}\hfill \\ & \times d\omega \frac{\mathrm{\Gamma }(\omega )n_{L/R}(\omega )}{\mathrm{\Delta }_{s_1^{}s_2}\mathrm{\Delta }B\omega +\text{i}0^+},\hfill \end{array}$$
(17)
where $`\mathrm{\Delta }_{s_1^{}s_2}=E_{s_1^{}}E_{s_2}`$ and $`E_s`$ are the eigenvalues of $`_{\mathrm{d}\mathrm{o}\mathrm{t}}`$. $`n_r(\omega )=1/\left(e^{\beta (\omega \mu _r)}1\right)`$ is the Bose function and the spectral function $`\mathrm{\Gamma }(\omega )`$ is given by
$$\mathrm{\Gamma }(\omega )=\pi S\underset{k}{}|1|k|^2\delta (\omega \omega _k),$$
(18)
with $`\omega _k=2JS[1\mathrm{cos}(ka)]+B_0`$. In (17) $`\mathrm{\Delta }B`$ corresponds to the left/right reservoir. Due to the conservation of the total spin the following equation holds for the $`s^z`$ quantum numbers of all non-vanishing diagrams, as indicated in Fig. 1(b):
$$m_{s_1}m_{s_1^{}}=m_{s_2}m_{s_2^{}}.$$
(19)
Assuming the dot to be in thermal equilibrium initially, the non-diagonal matrix elements $`p(0)_{s,s^{}}`$ of the initial reduced density matrix vanish. From (19) it then follows that the only non-trivial non-diagonal elements of the reduced density matrix $`p(t)`$ (particularly $`p^{\text{st}}`$) generated are those between the two states with $`m=0`$.
In some cases the influence of these non-diagonal elements onto the stationary spin current may be neglected. Then (15) and (16) reduce to classical master equations:
$$\underset{r,s^{}}{}\underset{q=\pm }{}\left(\mathrm{\Sigma }_{s,s^{}}^{rq}p_s^{}^{\text{st}}\mathrm{\Sigma }_{s^{},s}^{rq}p_s^{\text{st}}\right)=0,$$
(20)
$$I_r^{\text{st}}=\underset{s,s^{}}{}\left(\mathrm{\Sigma }_{s,s^{}}^{r+}p_s^{}^{\text{st}}\mathrm{\Sigma }_{s^{},s}^rp_s^{\text{st}}\right),$$
(21)
where $`q=\pm `$ corresponds to transitions with a spin-flip of $`\mathrm{\Delta }m=\pm 1`$ in the dot. The rates $`\mathrm{\Sigma }_{s,s^{}}^{r\pm }`$ are given by
$`\mathrm{\Sigma }_{s,s^{}}^{L/R+}`$ $`=\left|s|_iJ_{L/R}^is_i^+|s^{}\right|^2`$
$`\times \mathrm{\Gamma }(\mathrm{\Delta }_{ss^{}}\mathrm{\Delta }B)n_{L/R}(\mathrm{\Delta }_{ss^{}}\mathrm{\Delta }B),`$ (22a)
$`\mathrm{\Sigma }_{s,s^{}}^{L/R}`$ $`=\left|s^{}|_iJ_{L/R}^is_i^+|s\right|^2`$
$`\times \mathrm{\Gamma }(\mathrm{\Delta }_{s^{}s}\mathrm{\Delta }B)\left[1+n_{L/R}(\mathrm{\Delta }_{s^{}s}\mathrm{\Delta }B)\right],`$ (22b)
corresponding to absorption and emission of magnons by the dot, respectively. This is the general form for arbitrary $`\mathrm{\Delta }B`$ and $`\mu _r`$ (the latter enters the Bose function). For investigation of the two driving mechanisms we set either $`\mathrm{\Delta }B=0`$ ($`\mathrm{\Delta }\mu `$-configuration) or $`\mu _L=\mu _R=0`$ ($`\mathrm{\Delta }B`$-configuration). The product of the spectral- and Bose function, entering (22a), is depicted in Fig. 9 (left). While the spectral function introduces the band edges, the Bose function leads to the exponential decay of the rates.
## III Results
First we state some basic principles determining the single magnon transport through the dot, for the case of parallel magnetized reservoirs (the modifications for the case of antiparallel magnetized reservoirs will be stated in section III.3).
Tunneling of a magnon out of/into a reservoir involves a spin-flip $`\mathrm{\Delta }m=\pm 1`$ in the dot. Moreover it requires an transition energy above the lower band edges $`\omega _{L/R}`$ of the corresponding reservoirs.<sup>4</sup><sup>4</sup>4 For the parameters in consideration the upper band edge is of minor importance because there the magnon excitation is already suppressed exponentially. The magnon dispersion relation enters the rates via the spectral function $`\mathrm{\Gamma }(\omega )`$. At last the tunnel couplings $`J_r^i`$, which enter the rates as prefactors to the spin matrix elements, impose further restrictions onto the tunneling transitions. Fig. 2 shows the two coupling configurations studied, which we refer to as *serial* and *parallel coupling*.
For instance these may be realized by two spin-$`1/2`$ quantum dots coupled to the reservoirs in series or in parallel, respectively. Particularly for the case of parallel coupling we have examined the influence of a relative phase between the couplings. For vanishing phase no qualitatively new features occur in comparison to the serial coupling case. But for a phase of $`\pi `$ new features, such as the 2-magnon tunneling, arise. Therefore we restrict the discussion of the parallel coupling to this case, which is depicted in Fig. 2(b).
Depending on the couplings some states may decouple from either of the reservoirs. E.g., in the case of *parallel coupling* transitions between $`|𝖲`$ and $`|𝖳_{}`$ can occur over the left reservoir only, while transition between $`|𝖳_\mathrm{𝟢}`$ and $`|𝖳_{}`$ occur exclusively over the right one; see Fig. 2(b).
Fig. 3(a) shows the dot level structure, exemplarily for $`|𝖲`$/$`|𝖳_\mathrm{𝟢}`$-eigenstates, i.e. $`\mathrm{\Delta }H=0`$, and its dependence on the dot magnetic field $`B_{\text{Zee}}`$ and the intra-dot exchange coupling $`J_{\mathrm{d}\mathrm{o}\mathrm{t}}`$. The average $`\overline{H}=(H_1+H_2)/2`$ of the effective field $`H_i`$ adds to the homogeneous magnetic field $`B_{\text{Zee}}`$. The possible transitions in the case of serial coupling are illustrated in Fig. 3(b) \[note that transitions can only occur between states where the state with higher energy has also a higher spin quantum number in $`z`$direction, since absorption (emission) of a magnon by the dot increases (decreases) energy and spin simultaneously\]. In this case, the couplings impose no further restrictions, i.e., if the transition energy between the dot levels is sufficient, the respective transition can occur over both of the reservoirs.
In view of the spin current the various transitions contribute differently; see Fig. 4, where only the essential current contributions are shown. Due to the exponential decay of the magnon occupation with energy, large transition energies lead to small current contributions due to suppressed magnon absorption processes (note that current can only occur if there is a closed cycle of magnon absorption and emission processes). Therefore the $`𝖳_+𝖲`$-contribution is left out.
Furthermore the upper levels are generally less occupied and thus play also a minor role in current transport, as long as other ‘channels’ are available, as is the case in Fig. 4. Since the stationary dot spin is conserved it is sufficient to consider the current $`I_L`$ over the left reservoir. W.l.o.g. we choose $`\mu _{L/R}`$, resp. $`\mathrm{\Delta }B`$ such, that the spin current flows from the left to the right reservoir. Due to its insignificance the backflow from the right to the left is also not depicted in Fig. 4.
### III.1 Spin current for $`\mathrm{\Delta }\mu `$-configuration
Serial coupling. For the $`\mathrm{\Delta }\mu `$-configuration the band edges are given by $`\omega _L=\omega _R=0`$. Therefore transitions are possible for arbitrary small transition energies. Fig. 5 shows the spin current for the $`\mathrm{\Delta }\mu `$-configuration in serial coupling as a function of the dot magnetic field $`B_{\text{Zee}}`$ and the intra-dot exchange coupling $`J_{\mathrm{d}\mathrm{o}\mathrm{t}}`$.
With each new channel opened the spin current increases. For $`J_{\mathrm{d}\mathrm{o}\mathrm{t}}=0`$ the current vanishes, since in this case the whole system separates into two uncoupled parts (left and right); see Fig. 2(a). The cross section indicated by the arrow is shown in Fig. 6, where the level schemes illustrate the main current contributions. Each time the singlet level crosses one of the triplet levels, a qualitative change of the spin current characteristics occurs.
The outstanding current peak is due to the onset of the $`\mathrm{𝖲𝖳}_{}`$-transition, opening a new current channel which involves the strongly occupied $`𝖳_{}`$-level.
The width of the dip at $`J_{\mathrm{d}\mathrm{o}\mathrm{t}}=0`$ is determined by the inhomogeneity $`\mathrm{\Delta }H`$ of the effective field $`H_i`$; see upper inset in Fig. 6. To understand this, consider the $`J_{\mathrm{d}\mathrm{o}\mathrm{t}}=0`$ eigenstates $`|`$/$`|`$ of $`_{\mathrm{d}\mathrm{o}\mathrm{t}}`$, to which a finite $`J_{\mathrm{d}\mathrm{o}\mathrm{t}}\mathrm{\Delta }H`$ may be regarded as a perturbation. The corresponding level scheme is depicted in the lower inset. Here, each transition may exclusively occur over one of the reservoirs only. As a consequence the stationary current vanishes for $`J_{\mathrm{d}\mathrm{o}\mathrm{t}}=0`$, since there exists no path leading back to an arbitrary initial state while at the same time carrying spin from the left to the right reservoir. E.g., on the path $`|\stackrel{𝐿}{}|𝖳_+\stackrel{𝑅}{}|`$ a magnon tunnels from left to right. But in order to get back to the initial state, which is necessary to get a stationary current, a magnon must tunnel back again. So effectively the stationary spin current vanishes. Switching on $`J_{\mathrm{d}\mathrm{o}\mathrm{t}}`$ evokes a ‘coupling’ between the $`|`$/$`|`$-states. Since this coupling opens a new channel between $`|`$ and $`|`$ without magnon flow, the paths indicated in the inset level scheme give rise to a nonvanishing spin current. The $`J_{\mathrm{d}\mathrm{o}\mathrm{t}}`$-mediated coupling becomes important for $`J_{\mathrm{d}\mathrm{o}\mathrm{t}}`$ of the order of $`\mathrm{\Delta }H`$, which measures the splitting between the two levels. Thus $`\mathrm{\Delta }H`$ determines the width of the current dip.
For $`|\mathrm{\Delta }H||\mathrm{min}\{J_r^i\}|`$, particularly $`\mathrm{\Delta }H=0`$, a narrow dip still remains. In this case the eigenstates are $`𝖲`$/$`𝖳_\mathrm{𝟢}`$-like also in the neighborhood of $`J_{\mathrm{d}\mathrm{o}\mathrm{t}}=0`$. Then, the suppression of the spin current in the dip is due to a destructive interference between equivalent current carrying paths, e.g. the $`\mathrm{𝖲𝖳}_{}`$\- and $`𝖳_\mathrm{𝟢}𝖳_{}`$-paths. Thus, here the nondiagonal matrixelement $`p_{\mathrm{𝖲𝖳}_\mathrm{𝟢}}^{\text{st}}`$ has to be taken into account, i.e. only the full kinetic equation yields correct results in the neighborhood of $`J_{\mathrm{d}\mathrm{o}\mathrm{t}}=0`$.
Tuning further up to $`J_{\mathrm{d}\mathrm{o}\mathrm{t}}>B_{\text{Zee}}`$, the spin current finally saturates in the contribution of the $`𝖳_\mathrm{𝟢}𝖳_{}`$-transition, while the $`𝖲`$-level occupation gets negligible.
Parallel coupling. Fig. 7 shows the spin current for parallel coupling as a function of $`B_{\text{Zee}}`$ and $`J_{\mathrm{d}\mathrm{o}\mathrm{t}}`$.
For $`J_{\mathrm{d}\mathrm{o}\mathrm{t}}=0`$ instead of a dip, now a *peak* occurs in the spin current. In contrast to the current peak in Fig. 6 this peak cannot be ascribed to the onset of a transition. To understand this we discuss the current along the cross-section indicated; see Fig. 8.
Again, the level schemes illustrate the largest contributions to the spin current. The difference between serial and parallel coupling is, that in the case of the latter the transitions are more restricted \[as mentioned in the context of Fig. 2(b)\]. Transitions involving $`|𝖲`$ ($`|𝖳_\mathrm{𝟢}`$) may occur exclusively over the left (right) reservoir here. Therefore a stationary spin current requires a sequential tunneling of *two* magnons, e.g. on the path $`|𝖳_{}\stackrel{𝐿}{}|𝖲\stackrel{𝐿}{}|𝖳_+\stackrel{𝑅}{}|𝖳_\mathrm{𝟢}\stackrel{𝑅}{}|𝖳_{}`$, as illustrated in the level scheme for sector $`(d)`$.
The current peak $`(c)`$ appears for $`J_{\mathrm{d}\mathrm{o}\mathrm{t}}\mathrm{\Delta }H`$ since then the inhomogeneity $`\mathrm{\Delta }H`$ becomes essential by ‘coupling’ the $`|𝖲`$/$`|𝖳_\mathrm{𝟢}`$-states and thus introducing a further channel which enables 1-magnon tunneling; see the central level scheme. In this way the otherwise inevitable 2-magnon tunneling is cut short, resulting in an increased current \[compared to the 1-magnon tunneling peak $`(c)`$ the current in sector $`(d)`$ is exponentially decreased by a factor of $`\mathrm{exp}(\beta \mathrm{\Delta }_{𝖳_\mathrm{𝟢}𝖳_{}})`$\].
Unlike in sector $`(d)`$ the $`\mathrm{\Delta }H`$-mediated coupling still takes influence in sector $`(b)`$, as illustrated in the according level scheme. While in the former case the $`\mathrm{𝖲𝖳}_{}`$-transition is the ‘bottleneck’, which cannot be short-circuited, in the latter case the restricting $`𝖳_+𝖲`$-transition is cut short indeed.
Besides these qualitative arguments it shall be mentioned that the quantitative computation of the spin current is performed by setting up the kinetic equation in terms of the eigenbasis of $`_{\mathrm{d}\mathrm{o}\mathrm{t}}`$. The influence of nondiagonal matrixelements between the two states with $`m=0`$, which gets most important for a small level splitting, has to be considered too. For $`J_{\mathrm{d}\mathrm{o}\mathrm{t}}=0`$ the eigenstates are $`|`$/$`|`$ while their minimal splitting is given by $`\mathrm{\Delta }H`$. Since the inhomogeneity $`\mathrm{\Delta }H`$ cannot vanish for the $`J_r^i`$ in parallel coupling, the states are never degenerate and thus a destructive interference is largely suppressed. However a small current correction still remains, resulting in a reduced peak height. Thus, in order to account for this correction, the full kinetic equation has to be considered instead of the classical master equation.
### III.2 Spin current for $`\mathrm{\Delta }B`$-configuration
In the case of the $`\mathrm{\Delta }B`$-configuration a magnetic field gradient drives the spin current. Here the minimal excitation energies in the left and right reservoir differ, since the band edges are shifted by the magnetic field gradient: $`\omega _{L/R}=B_0\pm \mathrm{\Delta }B`$. Therefore, in contrast to the $`\mathrm{\Delta }\mu `$-configuration, the transitions over the left and right reservoir set in for *different* dot transition energies which leads to additional features in the spin current. The energy dependence of the rates entering the kinetic equation, particularly the shift of the band edges, is shown in Fig. 9. It also illustrates the successive onset of the transitions over the right and left reservoir, exemplarily for the $`\mathrm{𝖲𝖳}_{}`$-transition.
Serial coupling. Fig. 10 shows the spin current for the $`\mathrm{\Delta }B`$-configuration in serial coupling.
Since we choose $`\mathrm{\Delta }B=B_0`$ here, it follows that $`\omega _R=0`$. However, a different choice does not lead to relevant modifications. Again the onsets of transitions are indicated by lines. For $`J_{\mathrm{d}\mathrm{o}\mathrm{t}}\mathrm{\Delta }H`$ deviations from these lines occur due to the influence of the inhomogeneity $`\mathrm{\Delta }H`$, which shifts the eigenenergies of $`_{\mathrm{d}\mathrm{o}\mathrm{t}}`$ non-linearly.
In sector $`(\text{I})`$ the current is basically the same as for the case of the $`\mathrm{\Delta }\mu `$-configuration, compare Fig. 5. Beyond this similarity new features arise, specific to the $`\mathrm{\Delta }B`$-configuration. In the lower half of Fig. 10 an area of zero current occurs, denoted as sector $`(\text{II})`$. The left reservoir is completely decoupled here since all transitions involving it are disabled (all the transition energies lie below the left band edge). Furthermore there are additional steps in the sectors $`(\text{III})`$ and $`(\text{IV})`$. To understand these, consider Fig. 11 which shows the cross section indicated.
In sector $`(\text{III})`$ the onset of the $`\mathrm{𝖲𝖳}_{}`$-transition over the right reservoir \[compare Fig. 9(b)\] does not open a new current channel but leads to a swap of occupation between the $`𝖲`$\- and $`𝖳_{}`$-level; see inset. Thus the increase of the spin current is due to the enhanced $`𝖳_{}`$-occupation. This is revealed by the fact that the $`𝖳_{}`$-occupation is effectively mapped onto the current (compare inset). An analog argumentation holds for the step in sector $`(\text{IV})`$.
In contrast to this, the peak $`(\text{I\hspace{0.17em}a})`$ arises since a new (resonant) current channel is opened by the onset of the $`\mathrm{𝖲𝖳}_{}`$-transition over the left reservoir. This is completely analog to the corresponding peak in the case of the $`\mathrm{\Delta }\mu `$-configuration (compare Fig. 6). For sectors $`(\text{I\hspace{0.17em}b})`$ and $`(\text{I\hspace{0.17em}c})`$ there is no qualitative difference as well since all transitions are possible over both reservoirs, as is the case for the $`\mathrm{\Delta }\mu `$-configuration too.
Parallel coupling. Analog to the case of the $`\mathrm{\Delta }\mu `$-configuration (compare Fig. 5 and 7) for $`J_{\mathrm{d}\mathrm{o}\mathrm{t}}=0`$ the dip is replaced by a peak in sector $`(\text{I})`$. Nevertheless the sector $`(\text{II})`$ of zero current, as also the steps in sectors $`(\text{III})`$ and $`(\text{IV})`$ remain as for the case of serial coupling.
Finally, for the case of serial coupling, we discuss the spin current as a function of the *magnetic field gradient* $`\mathrm{\Delta }B`$; see Fig. 12, where the dot parameters $`B_{\text{Zee}}`$ and $`J_{\mathrm{d}\mathrm{o}\mathrm{t}}`$ are fixed.
For small $`\mathrm{\Delta }B`$ the current increases linearly. Then, due to the shift of the band edges, the transitions over the left reservoir set out successively (see inset), each resulting in a drop of the current. When the left reservoir finally decouples, the current vanishes completely.
The shape of the current as function of $`\mathrm{\Delta }B`$ depends strongly on the dot level structure determining which transitions are involved in which order. Fig. 13 shows exemplarily the current as a function of $`\mathrm{\Delta }B`$ and $`J_{\mathrm{d}\mathrm{o}\mathrm{t}}`$. Once again with each onset of a transition the current increases.
### III.3 Antiparallel magnetized reservoirs
So far we have considered the reservoirs to be magnetized in the same direction. In the following we briefly discuss the case of two antiparallel magnetized reservoirs, where w.l.o.g. the ground state of the right reservoir is flipped upwards. The real-time formalism used is also capable to deal with such a system. This is simply achieved by exchanging the magnon creation and annihilation operators for the right reservoir in (5), (6) and (11).<sup>5</sup><sup>5</sup>5This results in exchanging $`n_r^+n_r^{}`$ and reversing the sign of $`\omega \omega `$ in the resolvent of the expression (25) \[the latter involves a sign reversal of the arguments of the spectral- and Bose functions in (22)\].
In this antiparallel configuration a spin current flows even without an additional driving mechanism, such as a magnetic field gradient applied, since the reservoir magnetizations tend to equalize by exciting magnons in both of the reservoirs. Fig. 14 shows exemplarily the spin current as a function of $`J_{\mathrm{d}\mathrm{o}\mathrm{t}}`$ for the case of serial coupling. In view of the basic transport principles the main change is, that transitions over the right reservoir can now only occur between states where the state with higher energy has the lower spin quantum number in $`z`$direction, since the magnons in the right reservoir carry a ‘spin-down’.
The spin current $`I_L`$ is negative since ‘spin-ups’ are entering the left reservoir, in order to equalize the magnetizations. The current is suppressed, when the singlet level lies between the two outer triplet levels, since then the right reservoir is completely decoupled. In the remaining sectors a non-vanishing spin current is carried by magnon tunneling processes which involve all dot levels, as indicated in the level schemes in Fig. 14. Since these processes require a magnon to be initially excited in the left reservoir (long upward arrow in the level schemes), the current decays exponentially with increasing $`|J_{\mathrm{d}\mathrm{o}\mathrm{t}}|`$. Furthermore, due to this the current vanishes with vanishing temperature.
We conclude that also in this antiparallel configuration the spin current can be switched by tuning the intra-dot parameters $`B_{\text{Zee}}`$ and $`J_{\mathrm{d}\mathrm{o}\mathrm{t}}`$, e.g. by simply evoking a $`\mathrm{𝖲𝖳}_{}`$-crossing.
### Conclusions
In summary, our investigation of the proposed spin quantum dot system revealed a rich structure of the spin current as a function of the intra-dot parameters $`B_{\text{Zee}}`$ and $`J_{\mathrm{d}\mathrm{o}\mathrm{t}}`$. Due to the involved large current variations, a switching of the current is possible in principle.
For the case of *parallel magnetized reservoirs* the influence of the dot-reservoir couplings on the spin transport becomes apparent by opposing the serial and parallel coupling configurations. The restriction of the transitions in the case of parallel coupling leads to drastically different transport characteristics, which originate from the sequential 2-magnon tunneling. However, due to the differing signs of the dot-reservoir couplings, the parallel coupling configuration represents rather an extreme case compared to the serial coupling configuration which suggests itself. It shall be mentioned that small deviations from the ‘pure’ coupling configurations discussed here, do not lead to qualitative modifications of the spin current results.
The concept of a magnetochemical potential as sum of the chemical potential and the magnetic field, analog to the electrochemical potential in the electronic case, is deficient. The magnon chemical potential controls the magnon occupation, whereas an applied magnetic field gradient shifts the magnon band edges while leaving the overall magnon occupation unchanged. As we have shown, both mechanisms, the differing left and right magnon chemical potentials as also the magnetic field gradient, yield different results for the spin current. This is due to the fact that the magnon dispersion relation, particularly the band edges, play an essential role in spin transport, in contrast to the case of electronic transport, which primarily takes place between the electrochemical potentials.
In the case of *antiparallel magnetized reservoirs* the main features of the spin current can be explained by considering slightly modified spin transport principles. As we have shown, here spin current switching is possible too.
We remark, that in those sectors where the spin current carried by single magnon tunneling is exponentially suppressed, a cotunneling of magnons, analog to its electronic pendant, may yield non-negligible contributions to the spin current. However, the current suppression due to magnon dispersion effects should remain untouched.
We gratefully acknowledge discussions with F. Meier, J. Martinek, D. Loss and H. Capellmann.
*
## Appendix A
The diagrams contributing to the kernels $`\mathrm{\Sigma }(\text{i}0^+)`$ and $`\mathrm{\Sigma }_{I_r}(\text{i}0^+)`$ to lowest order in tunneling are shown in Fig. 15.
Thereby, in matrix representation the kernels are given by
$`\mathrm{\Sigma }(\text{i}0^+)_{s_1s_1^{},s_2s_2^{}}`$ $`={\displaystyle \underset{r=L,R}{}}{\displaystyle \underset{d=1}{\overset{8}{}}}d`$ (23)
$`\mathrm{\Sigma }_{I_r}(\text{i}0^+)_{s_1s_1^{},s_2s_2^{}}`$ $`=1+2+5+6.`$ (24)
By defining the auxiliary function
$$\sigma _{ij}^{rp}(\mathrm{\Delta }E)=\frac{\text{i}}{2\pi }𝑑\omega \frac{\mathrm{\Gamma }(\omega )n_r^p(\omega )}{\mathrm{\Delta }E\overline{B}_r\omega +\text{i}0^+}$$
(25)
where $`p\{+,\}`$, $`n_r^+(\omega )=n_r(\omega )`$, $`n_r^{}(\omega )=1+n_r(\omega )`$ and $`\overline{B}_{L/R}=\pm \mathrm{\Delta }B`$, the expressions corresponding to the diagrams with left running reservoir contractions are given as follows:
$`1`$ $`={\displaystyle \underset{i,j}{}}s_1|J_r^is_i^+|s_2s_2^{}|J_r^js_j^{}|s_1^{}\sigma _{ij}^{r+}(\mathrm{\Delta }_{s_1^{}s_2}),`$
$`3`$ $`={\displaystyle \underset{i,j}{}}s_1|J_r^js_j^{}|s_2s_2^{}|J_r^is_i^+|s_1^{}\sigma _{ij}^r(\mathrm{\Delta }_{s_2^{}s_1}),`$
$`5`$ $`={\displaystyle \underset{i,j}{}}{\displaystyle \underset{s}{}}\delta _{s_1^{}s_2^{}}s_1|J_r^is_i^+|ss|J_r^js_j^{}|s_2\sigma _{ij}^r(\mathrm{\Delta }_{s_1^{}s}),`$
$`7`$ $`={\displaystyle \underset{i,j}{}}{\displaystyle \underset{s}{}}\delta _{s_1s_2}s_2^{}|J_r^js_j^{}|ss|J_r^is_i^+|s_1^{}\sigma _{ij}^{r+}(\mathrm{\Delta }_{ss_1}).`$
Each of the remaining diagrams $`d`$ with right running contraction is obtained by taking the complex conjugate and interchanging the primed by unprimed indices (and vice versa) in the expression for the left running diagram $`d1`$. E.g.:
$$2=\left[\underset{i,j}{}s_1^{}|J_r^is_i^+|s_2^{}s_2|J_r^js_j^{}|s_1\sigma _{ij}^{r+}(\mathrm{\Delta }_{s_1s_2^{}})\right]^{}.$$
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