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# Enumeration of ramified coverings of the sphere and 2-dimensional gravity
## 1 Introduction
Denote by $`๐`$ the subalgebra of the algebra of power series in one variable, generated by the series
$$\underset{n1}{}\frac{n^{n1}}{n!}q^n\text{and}\underset{n1}{}\frac{n^n}{n!}q^n.$$
We wish to show that this algebra plays an important role in the intersection theory of moduli spaces $`\overline{}_{g,n}`$ of stable curves and in the problem of enumeration of ramified coverings of the sphere.
#### SECTION 2
contains a more explicit description of the algebra $`๐`$. In this section we also prove some relations between $`๐`$ and the combinatorics of Cayley trees ($`=`$ trees with numbered vertices).
#### SECTION 3
is devoted to the problem of enumerating the ramified coverings of the sphere with specified ramification types.
Consider a holomorphic map $`f:C\mathrm{P}^1`$ of degree $`n`$ from a smooth complex curve $`C`$ to the Riemann sphere. Such maps will be called ramified coverings with $`n`$ sheets.
A ramification point of $`f`$ is a point of the target Riemann sphere that has less than $`n`$ distinct preimages.
For each ramification point $`y`$ of a ramified covering, we are going to single out several simple preimages of $`y`$.
###### Definition 1.1
A marked ramified covering is a ramified covering with a choice, for every ramification point $`y`$, of a subset of the set of simple preimages of $`y`$.
Consider a partition $`\mu =1^{a_1}2^{a_2}\mathrm{}`$ of an integer $`mn`$. (Here we use multiplicative notation for partitions: the partition $`\mu `$ contains $`a_1`$ parts equal to $`1`$, $`a_2`$ parts equal to $`2`$, and so on, $`ia_i=m`$.) Suppose that a point $`y\mathrm{P}^1`$ has $`a_1`$ marked simple preimages, $`a_2`$ double preimages, and so on. (Consequently, $`y`$ also has $`nm`$ unmarked simple preimages.) We then say that $`y`$ is a ramification point of $`f`$ of multiplicity $`r=a_2+2a_3+3a_4+\mathrm{}`$ and of ramification type $`\mu =1^{a_1}2^{a_2}\mathrm{}`$. Sometimes the number $`r`$ will also be called the degeneracy of the partition $`\mu `$.
###### Definition 1.2
A Hurwitz number $`h_{n;\mu _1,\mathrm{},\mu _k}`$ is the number of connected $`n`$-sheeted marked ramified coverings of $`\mathrm{P}^1`$ with $`k`$ ramification points, whose ramification types are $`\mu _1,\mathrm{},\mu _k`$. Every such covering is counted with weight $`1/|\mathrm{Aut}|`$, where $`|\mathrm{Aut}|`$ is the number of automorphisms of the covering.
Note that the genus $`g`$ of the covering surface can be reconstituted from the data $`(n;\mu _1,\mathrm{},\mu _k)`$ using the Riemann-Hurwitz formula: if the degeneracy of $`\mu _i`$ equals $`r_i`$, then
$$22g=2nr_1\mathrm{}r_k.$$
Fix $`k`$ nonempty partitions $`\mu _1,\mathrm{},\mu _k`$ with degeneracies $`r_1,\mathrm{},r_k`$. Let $`r`$ be the sum $`r=r_1+\mathrm{}+r_k`$.
###### Notation 1.3
Denote by $`h_{g,n;\mu _1,\mathrm{},\mu _k}`$ the number of $`n`$-sheeted marked ramified coverings of $`\mathrm{P}^1`$ by a genus $`g`$ surface, with $`k`$ ramification points of types $`\mu _1,\mathrm{},\mu _k`$ and, in addition, $`c(n)=2n+2g2r`$ simple ($`=`$ of multiplicity 1) ramification points. Each covering is counted with weight $`1/|\mathrm{Aut}|`$.
###### Theorem 1
Fix any $`g0`$, $`k0`$. If $`g=1`$, we suppose that $`k1`$. Then for any partitions $`\mu _1,\mathrm{},\mu _k`$, the series
$$H_{g;\mu _1,\mathrm{},\mu _k}(q)=\underset{n1}{}\frac{h_{g,n;\mu _1,\mathrm{},\mu _k}}{c(n)!}q^n$$
lies in the algebra $`๐`$.
The only known proof of this theorem involves a surprising detour by the intersection theory on moduli spaces of stable curves. Using the Ekedahl-Lando-Shapiro-Vainshtein (or โELSVโ) formula, one can express the Hurwitz numbers for $`k=1`$ as integrals of some cohomology classes over these moduli spaces. In the ELSV formula is used to express the generating functions for Hurwitz numbers in the case $`k=1`$ as rational functions of the series $`Y(q)`$. This essentially proves the theorem for $`k=1`$. After that, we proceed by induction on $`k`$.
Among other things, the theorem allows one to find the asymptotic of the coefficients of $`H_{g;\mu _1,\mathrm{},\mu _k}`$ as $`n\mathrm{}`$, knowing only the several first coefficients.
#### SECTION 4
describes a model of $`2`$-dimensional gravity obtained by counting ramified coverings of the sphere. We show that one can extract from the symptotic of Hurwitz numbers a solution of the Painlevรฉ I equation and the string solution of the Korteweg - de Vries equation. The same solutions are obtained in other models of $`2`$-dimensional gravity (by counting quadrangulations or using intagrals over moduli spaces of curves).
We also compare the enumerative problems concerning ramified coverings of the sphere and those of the torus. While the former are related to the intersection theory on $`\overline{}_{g,n}`$ and give rise to the algebra $`๐`$, the latter are related to volumes of spaces of abelian differentials on Riemann surfaces and give rise to the algebra of quasi-modular forms.
#### Acknowledgments
The author is grateful to J.-M. Bismut, F. Labourie, M. Kontsevich, S. Natanzon, Ch. Okonek, A. Okounkov, D. Panov, J.-Y. Welschinger, D. Zagier, A. Zorich and A. Zvonkin for useful discussions and remarks. A special thank to Sergei Lando, with whom we proved together some of the results of the last section, and to M. Kazarian for sharing his own work on the same subject. I would also like to thank for their interest the participants of the mathematical physics seminar at the ETH Zรผrich and of the mathematical seminar at the ENS Lyon, as well as the participants of the Luminy conference on billiards and Teichmรผller spaces.
This work was partially supported by EAGER - European Algebraic Geometry Research Training Network, contract No. HPRN-CT-2000-00099 (BBW) and by the RFBR grant 02-01-22004.
#### Notation.
Here we summarize some notation that we use consistently throughout the paper.
| $`n`$ | The number of sheets of a covering. The power of the variable $`q`$ in generating series. The number of marked points on a Riemann surface is sometimes $`n`$ and sometimes $`nr`$. |
| --- | --- |
| $`q`$ | The variable in generating series (to a sequence $`s_n`$ we usually assign the series $`s_nq^n/n!`$). |
| $`g`$ | The genus of a Riemann surface. |
| $`\mu `$ | A partition. |
| $`p`$ | The number of parts of a partition $`\mu `$. |
| $`a_i`$ | The number of parts of a partition $`\mu `$ that are equal to $`i`$. |
| $`b_i`$ | The parts of a partition $`\mu `$ are denoted by $`b_1,\mathrm{},b_p`$. |
| $`r`$ | The degeneracy of a partition $`\mu `$ defined by $`r=(b_i1)`$. The multiplicity of a ramification point. |
| $`k`$ | The number of partitions. If $`k>1`$, the partitions are denoted by $`\mu _1,\mathrm{},\mu _k`$, their degeneracies by $`r_1,\mathrm{},r_k`$, while $`r`$ is the total degeneracy $`r=r_i`$. |
| $`c(n)`$ | The number of simple ramification points in a ramified covering. |
| $`\psi _i`$ | The first Chern class $`c_1(_i)`$ of the line bundle $`_i`$ over $`\overline{}_{g,n}`$. |
| $`d_i`$ | The power of the class $`\psi _i`$ in the intersection numbers we consider. |
## 2 The algebra $`๐`$ of power series
The algebra of power series
$$๐=[\underset{n1}{}\frac{n^{n1}}{n!}q^n,\underset{n1}{}\frac{n^n}{n!}q^n]$$
plays a central role in this paper. Here we give an explicit description of $`๐`$ and show its relation with the combinatorics of Cayley trees. Many of the results below are known, but have probably never been put together. As far as we know, the algebra $`๐`$ itself was first discovered by D. Zagier several years ago (unpublished), and then independently introduced in our paper , where most of the results of Section 2.1 are given. Various series from $`๐`$ also appear in .
### 2.1 How to make computations in $`๐`$
Denote by $`Y`$ and $`Z`$ the generators of $`๐`$
$$Y=\underset{n1}{}\frac{n^{n1}}{n!}q^n,Z=\underset{n1}{}\frac{n^n}{n!}q^n.$$
Denote by $`D`$ the differential operator $`D=q\frac{}{q}`$. Thus $`Z=DY`$.
Note that both $`Y`$ and $`Z`$ have a radius of convergence of $`1/e`$. Therefore the same is true of all series in $`๐`$. The function $`Y(q)`$, more precisely, $`Y(q)`$, was considered by J. H. Lambert in 1758<sup>1</sup><sup>1</sup>1We thank N. AโCampo for this reference.. The relations that follow can be deduced from the Lagrange inversion theorem applied to the equation $`Y(q)=qe^{Y(q)}`$ or from the Abel identities (see , Section 1.2).
###### Proposition 2.1
We have
$$Y=qe^Y.$$
#### Proof.
$`Y`$ is the exponential generating series for rooted Cayley trees (Definition 2.9). Therefore $`e^Y`$ is the exponential generating series for forests of rooted Cayley trees. Add a new vertex $``$ to such a forest and join $``$ to the root of each tree. We obtain a Cayley tree with root $``$. This operation is a one-to-one correspondence, hence $`Y=qe^Y`$. โ
###### Corollary 2.2
On the disc $`|q|<1/e`$, the function $`Y(q)`$ is the inverse of the function $`q(Y)=Y/e^Y`$.
###### Proposition 2.3
We have $`(1Y)(1+Z)=1`$.
#### Proof.
$$Z=DY=D(qe^Y)=qe^Y+qe^YDY=qe^Y(1+Z)=Y(1+Z).$$
Hence $`(1Y)(1+Z)=1`$. โ
###### Corollary 2.4
As an abstract algebra, $`๐`$ is isomorphic to $`[X,X^1]`$, where $`X=1Y`$.
###### Proposition 2.5
We have
$$Y^k=k\underset{n1}{}\frac{(n1)\mathrm{}(nk+1)n^{nk}}{n!}q^n=k\underset{nk}{}\frac{n^{nk1}}{(nk)!}q^n.$$
#### Proof.
Induction on $`k`$. For $`k=1`$ the assertion is true. To go from $`k`$ to $`k+1`$, one uses the equality
$$D\left(\frac{Y^{k+1}}{k+1}\frac{Y^k}{k}\right)=(Y^kY^{k1})DY=Y^{k1}(Y1)Z=Y^k.$$
It is compatible with our expressions for $`Y^k`$ and $`Y^{k+1}`$, which determines $`Y^{k+1}`$ up to a constant. But yhe constant term of $`Y^{k+1}`$ vanishes. โ
Now we study the powers of $`Z`$.
###### Definition 2.6
Denote by $`A_n`$ the sequence of integers
$$A_n=\underset{\stackrel{p+q=n}{p,q1}}{}\frac{n!}{p!q!}p^pq^q.$$
Its first terms are $`0,2,24,312,4720,\mathrm{}`$. We have
$$Z^2=\underset{n1}{}\frac{A_n}{n!}q^n.$$
One can show that
$$A_n=n!\underset{k=0}{\overset{n2}{}}\frac{n^k}{k!}\sqrt{\pi /2}n^{n+\frac{1}{2}}.$$
As far as we know, there is no simple expression for the powers of $`Z`$. However, we can prove that they are linear combinations of the series
$$D^kZ=\underset{n1}{}\frac{n^{n+k}}{n!}q^n\text{and}D^k(Z^2)=\underset{n1}{}\frac{n^kA_n}{n!}q^n.$$
###### Proposition 2.7
For any integer $`k0`$, the power series $`D^kZ`$ and $`D^k(Z^2)`$ are polynomials in $`Z`$ with positive integer coefficients, of degrees $`2k+1`$ and $`2k+2`$ respectively:
$`D^k(Z)`$ $`=`$ $`(2k1)!!Z^{2k+1}+\text{lower order terms},`$
$`D^k(Z^2)`$ $`=`$ $`(2k)!!Z^{2k+2}+\text{lower order terms}.`$
#### Proof.
Applying $`D`$ to both sides of the equality $`(1Y)(1+Z)=1`$ we get
$$Z(1+Z)+(1Y)DZ=0.$$
Thus
$$DZ=\frac{Z(1+Z)}{1Y}=Z(1+Z)^2.$$
Hence
$$D(Z^2)=2Z^2(1+Z)^2.$$
Now we proceed by induction on $`k`$. โ
###### Corollary 2.8
For any positive integer $`k`$, the power series $`Z^k`$ is a linear combination with rational coefficients of the first $`k`$ series from the list $`Z,Z^2,DZ,D(Z^2),D^2Z,D^2(Z^2),\mathrm{}`$.
From Proposition 2.5 and Corollary 2.8 we deduce the following theorem.
###### Theorem 2
The algebra $`๐`$ is spanned over $``$ by the power series
$$1,\underset{n1}{}\frac{n^{n+k}}{n!}q^n,k,\underset{n1}{}\frac{n^kA_n}{n!}q^n,k.$$
Note that the Stirling formula together with the asymptotic for the sequence $`A_n`$ allows one to determine the leading term of the asymptotic for the coefficients of any series in $`๐`$. We have
$$\frac{n^n}{n!}\frac{1}{\sqrt{2\pi n}}e^n,\frac{A_n}{n!}\frac{1}{2}e^n.$$
Note also that if, for some series $`F๐`$, we know in advance its degree in $`Y`$ and in $`Z`$, then we can reconstitute the series $`F`$ using only a finite number of its initial terms โ a very useful property for computer experiments.
Combining both remarks, we see that initial terms of the sequence of coefficients of $`F`$ determine the asymptotic of the sequence.
### 2.2 Dendrology
###### Definition 2.9
A Cayley tree is a tree with numbered vertices.
It is well-known (Cayley theorem) that there are $`n^{n2}`$ Cayley trees with $`n`$ vertices. Note that the corresponding exponential generating function
$$\underset{n1}{}\frac{n^{n2}}{n!}q^n$$
lies in the algebra $`๐`$.
Consider a Cayley tree $`T`$ with two marked vertices $`a`$ and $`b`$. Denote by $`l(T)`$ the distance between these vertices, i.e., the number of edges in the shortest path joining them.
###### Definition 2.10
Denote by $`m_{n,k}`$ and $`p_{n,k}`$ the sums
$$m_{n,k}=\underset{T}{}l(T)^k,p_{n,k}=\underset{T}{}\frac{l(T)(l(T)1)\mathrm{}(l(T)k+1)}{k!}$$
where the sum is taken over all Cayley trees $`T`$ with $`n`$ vertices, two of which are marked.
For instance, $`m_{2,1}=p_{2,1}=2`$. Note that if we consider $`l(T)`$ as a random variable, then $`m_{n,k}`$ is its $`k`$th moment.
###### Theorem 3
For any $`k`$, the power series
$$\underset{n1}{}\frac{m_{n,k}}{n!}q^n\text{and}\underset{n1}{}\frac{p_{n,k}}{n!}q^n$$
lie in $`๐`$.
###### Example 2.11
It follows from the proof below that $`p_{n,1}=m_{n,1}=A_n`$. This number is called the total height of Cayley trees and was introduced in .
#### Proof of Theorem 3.
It is sufficient to prove the theorem for $`p_{n,k}`$.
Fix $`k`$. There is a natural bijection between the following sets of objects.
$`E_n`$ is the set of Cayley trees with $`n`$ vertices, on which one has marked two vertices by $`a`$ and $`b`$ and chosen $`k`$ distinct edges on the shortest path from $`a`$ to $`b`$. The number of elements in $`E_n`$ equals $`p_{n,k}`$.
$`F_n`$ is the set of ordered $`(k+1)`$-tuples of trees with $`n`$ vertices in whole; the vertices are numbered from $`1`$ to $`n`$ and, in addition, two vertices $`a_i`$ and $`b_i`$, $`1ik+1`$, are marked on each tree.
The bijection is established as follows. Take a forest from the set $`F_n`$. Draw new edges $`(b_1,a_2)`$, $`(b_2,a_3)`$, โฆ, $`(b_k,a_{k+1})`$. We obtain a tree with $`k`$ marked edges lying on the path between $`a_1`$ and $`b_{k+1}`$, i.e., a tree from the set $`E_n`$.
Now, the trees with two marked vertices are enumerated by the series $`Z`$, therefore the exponential generating series for the sequence $`|F_n|`$ is $`Z^{k+1}`$. โ
## 3 Counting ramified coverings of the sphere
This section is devoted to the enumeration of ramified coverings of the sphere by surfaces of a fixed genus $`g`$ and to a proof of Theorem 1.
### 3.1 The ELSV formula
Curiously, the most difficult part of the proof of Theorem 1 is the case with only one multiple ramification point, $`k=1`$. We know no other way to prove it than to use the intersection theory on moduli spaces. The main ingredient of the proof is a theorem by T. Ekedahl, S. K. Lando, M. Shapiro, and A. Vainshtein that we formulate below after introducing some notation.
Let $`\mu =1^{a_1}2^{a_2}\mathrm{}`$ be a partition with degeneracy $`r`$. We define $`|\mathrm{Aut}(\mu )|`$ to be $`|\mathrm{Aut}(\mu )|=a_1!a_2!\mathrm{}`$. For the formulation of the theorem it is more convenient to switch to using the additive notation for the partition $`\mu `$, $`\mu =(b_1,\mathrm{},b_p)`$, the $`b_i`$ being the parts of $`\mu `$. The Hurwitz number $`h_{g,n;\mu }`$ is defined in Notation 1.3.
We denote by $`_{g,n}`$ the moduli space of smooth genus $`g`$ curves with $`n`$ marked and numbered distinct points.
Further, $`\overline{}_{g,n}`$ is the Deligne-Mumford compactification of this moduli space; in other words, $`\overline{}_{g,n}`$ is the space of stable genus $`g`$ curves with $`n`$ marked points.
We denote by $`_i`$, $`1in`$, the $`i`$th tautological line bundle over $`\overline{}_{g,n}`$: consider a point $`x\overline{}_{g,n}`$ and the corresponding stable curve $`C_x`$; then the fiber of $`_i`$ over $`x`$ is the cotangent line to $`C_x`$ at the $`i`$th marked point. The first Chern class of $`_i`$ is denoted by $`c_1(_i)=\psi _i`$.
We will use the expression
$$\frac{1}{1\psi _i}=1+\psi _i+\psi _i^2+\mathrm{}H^{}(\overline{}_{g,n},).$$
Further, we introduce the Hodge vector bundle $`W`$ over $`\overline{}_{g,nr}`$. The fiber of $`W`$ over a smooth curve is the set of holomorphic $`1`$-forms on this curve. The fiber of $`W`$ over a general stable curve is the set of global sections of its dualizing sheaf. We do not give the details here (see the paper itself). Suffice it to note that $`W`$ is a vector bundle of rank $`g`$.
Now we can write down the ELSV formula.
###### Theorem 4
(The ELSV formula, ) For any $`g`$, $`n`$, and $`\mu `$ such that $`22g(nr)<0`$, we have
$$h_{g,n;\mu }=\frac{(2n+2g2r)!}{|\mathrm{Aut}(\mu )|}\underset{i=1}{\overset{p}{}}\frac{b_i^{b_i}}{b_i!}\times $$
$$\times \frac{1}{(npr)!}_{\overline{}_{g,nr}}\frac{c(W^{})}{(1b_1\psi _1)\mathrm{}(1b_p\psi _p)(1\psi _{p+1})\mathrm{}(1\psi _{nr})}.$$
### 3.2 Proof of Theorem 1
We are going to prove that all generating series $`H_{g;\mu _1,\mathrm{},\mu _k}`$ for the Hurwitz numbers $`h_{g,n;\mu _1,\mathrm{},\mu _k}`$ (see Notation 1.3) with respect to the number of sheets $`n`$ lie, once again, in the algebra $`๐`$. In Section 4 we give a motivation for considering these particular generating series.
First consider the case of just one partition $`k=1`$. This case was essentially covered in . Recently, M. Kazarian suggested an improvement of the theorem for $`k=1`$, giving an explicit expression for the series $`H_{g;\mu }`$ in terms of the generators $`Y`$ and $`Z`$ of $`๐`$.
###### Theorem 5
\[Kazarian\] Consider a partition $`\mu =(b_1,\mathrm{},b_p)`$, and let $`m=b_i`$. Then we have
$$H_{g;\mu }=\frac{1}{|\mathrm{Aut}(\mu )|}\underset{i=1}{\overset{p}{}}\frac{b_i^{b_i}}{b_i!}Y^m(Z+1)^{2g2+p}\phi (Z),$$
where $`\phi (Z)`$ is the polynomial
$$\phi (Z)=\underset{l0}{}\frac{Z^l}{l!}\underset{\overline{}_{g,p+l}}{}\frac{c(W^{})}{(1b_1\psi _1)\mathrm{}(1b_p\psi _p)}\frac{\psi _{p+1}^2\mathrm{}\psi _{p+l}^2}{(1\psi _{p+1})\mathrm{}(1\psi _{p+l})}.$$
Note that the last sum only goes up to $`l=3g3+p`$, otherwise the intergral equals $`0`$ for dimension reasons.
#### Sketch of a proof
(borrowed from ). Introduce the series $`F=F_{g;b_1,\mathrm{},b_p}`$ in an infinite number of variables
$$F(t_0,t_1,\mathrm{})=\underset{l;d_1,\mathrm{},d_l}{}\frac{t_{d_1}\mathrm{}t_{d_l}}{l!}\underset{\overline{}_{g,p+l}}{}\frac{c(W^{})\psi _{p+1}^{d_1}\mathrm{}\psi _{p+l}^{d_l}}{(1b_1\psi _1)\mathrm{}(1b_p\psi _p)}.$$
We have $`H_{g;b_1,\mathrm{},b_p}(q)=q^mF(q,q,q,\mathrm{})`$, where $`m=b_i`$. On the other hand, $`F`$ satisfies the string and the dilaton equations (see ):
$`{\displaystyle \frac{F}{t_0}}`$ $`=`$ $`mF+{\displaystyle \underset{i0}{}}t_{i+1}{\displaystyle \frac{F}{t_i}},`$
$`{\displaystyle \frac{F}{t_1}}`$ $`=`$ $`\chi F+{\displaystyle \underset{i0}{}}t_i{\displaystyle \frac{F}{t_i}},\chi =2g2+p.`$
The theorem now follows from the following fact, obtained by a manipulation of PDEs. For any series $`F`$ satisfying the above string and dilaton equations, let $`\phi (q)=F(0,0,q,q,q,\mathrm{})`$. Then we have
$$q^mF(q,q,q,\mathrm{})=Y^m(1+Z)^\chi \phi (Z).$$
Now we deduce the general case from the case $`k=1`$.
#### Theorem 2
Fix any $`g0`$, $`k0`$. If $`g=1`$, we suppose that $`k1`$. Then for any partitions $`\mu _1,\mathrm{},\mu _k`$, the series
$$H_{g;\mu _1,\mathrm{},\mu _k}(q)=\underset{n1}{}\frac{h_{g,n;\mu _1,\mathrm{},\mu _k}}{c(n)!}q^n$$
lies in the algebra $`๐`$.
#### Proof of Theorem 2.
The theorem is proved by induction on the number $`k`$ of partitions.
Base of induction. For $`k=0,1`$, the result is obtained by a direct application of Theorem 5. In the case $`k=0`$, we must use Theorem 5 with an empty partition $`\mu `$.
There are three exceptional cases in which Theorem 5 cannot be applied: $`g=0`$, $`k=0`$; $`g=0`$, $`k=1`$, $`p2`$; $`g=1`$, $`k=0`$. These cases are discussed in Remark 3.1 below. It turns out that the assertion of Theorem 1 fails only if $`g=1`$, $`k=0`$, as stated in the formulation.
Step of induction. The step of induction is an almost exact repetition of the proof of Theorem 2 from our previous work . We only give a short summary of the argument here. The proof goes in the spirit of . A similar proof, using the formalizm of colored permutations, is given in .
It is easy to see that there is only a finite number of possible cycle structures for a permutation that can be obtained as a product of two permutations with given cycle structures $`\mu _1`$ and $`\mu _2`$.
Let $`\mu _1`$ and $`\mu _2`$ be two partitions from the list $`\mu _1,\mathrm{},\mu _k`$. We can move the two corresponding ramification points on $`\mathrm{P}^1`$ towards each other until they collapse. We obtain a new (not necessarily connected) ramified covering. Its monodromy at the new ramification point is the product of the monodromies of the two points that have collapsed.
Let us choose one of the possible cycle structures of the product monodromy and also one of the possible ways in which the covering can split into connected components. By the induction assumption, we obtain a series from the algebra $`๐`$ assigned to each connected component of the covering. Indeed, each connected component is itself a ramified covering of the sphere as in Theorem 1, but with $`k1`$ fixed ramification types instead of $`k`$. We obtain the generating series for the number of nonconnected ramified coverings by multiplying the series that correspond to the connected components. Since it is a finite product of series lying in $`๐`$, we obtain again a series from $`๐`$.
Finally, we must add the generating series described above for all choices of types of nonconnected coverings. Since the number of choices is finite, we obtain, once again, a series from $`๐`$. โ
###### Remark 3.1
Let us consider the exceptional cases $`g=0`$, $`k=0,1`$ and $`g=1`$, $`k=0`$.
In the genus zero case, the ELSV formula transforms into a much simpler Hurwitz formula , which turns out to be applicable even if the multiple ramification point has only $`1`$ or $`2`$ preimages.
We have, using the notation of Theorem 4 and Notation 1.3,
$$h_{0,n;\mu }=\frac{(2n2r)!}{|\mathrm{Aut}(\mu )|}\underset{i=1}{\overset{p}{}}\frac{b_i^{b_i}}{b_i!}\frac{n^{nr3}}{(npr)!}.$$
This formula is true for any $`np+r`$ and for any partition $`\mu `$ (including even the empty partition). We see that the corresponding generating series always lies in the algebra $`๐`$.
The case $`g=1`$, $`k=0`$ is covered by the ELSV formula with an empty partition $`\mu `$. Consider the moduli space $`\overline{}_{1,1}`$. Denote by $`\beta `$ the $`2`$-cohomology class of $`\overline{}_{1,1}`$ whose integral over the fundamental homology class equals 1. One can prove that the Hodge bundle over $`\overline{}_{1,1}`$ is a line bundle with first Chern class $`\beta /24`$. Therefore we obtain
$$h_{1,n;\mathrm{}}=(2n)!\frac{1}{n!}_{\overline{}_{1,n}}\frac{1\frac{1}{24}\beta }{(1\psi _1)\mathrm{}(1\psi _n)}.$$
From this we get
$$\underset{n1}{}\frac{h_{1,n;\mathrm{}}}{(2n)!}q^n=\frac{1}{24}\underset{n1}{}\frac{A_n}{n}\frac{q^n}{n!}.$$
This series does not lie in $`๐`$ (and constitutes the only exception to the general rule). It suffices to consider the partition $`\mu =(1)`$, which amounts to distinguishing one sheet in the ramified covering, to obtain the series
$$\frac{1}{24}\underset{n1}{}\frac{A_n}{n!}q^n๐.$$
## 4 Random metrics and 2-dimensional gravity
In this section we propose a model of $`2`$-dimensional gravity via the enumeration of ramified coverings. We show that the โfree energyโ function and the values of โobservablesโ coinside with those obtained in other models.
We also draw a parallel between the study of spaces of Riemannian metrics using ramified coverings of the sphere and the study of spaces of abelian differentials using ramified coverings of the torus.
### 4.1 Models of 2-dimensional gravity
Here we explain what sort of questions about ramified coverings arise in $`2`$-dimensional gravity and why. Precise mathematical results are given below in Sections 4.2 and 4.3.
In every problem of statistical physics one starts with introducing a space of states and by assigning an energy to every state.
In 2-dimensional gravity, a state is a 2-dimensional compact oriented real not necessarily connected surface endowed with a Riemannian metric. Two surfaces like that are equivalent, i.e., correspond to the same state, if they are isometric.
Consider a surface $`S`$ with a Riemannian metric. Let $`\chi (S)`$ be its Euler characteristic and $`A`$ its total area. To such a surface one assigns an energy
$$E=\lambda A+\mu \chi (S).$$
Here $`\lambda `$ and $`\mu `$ are two constants called the cosmological constant and the gravitational constant, respectively. Note that $`\chi (S)`$ is actually the integral over $`S`$ of the scalar curvature of the metric. The fact that this integral takes such a simple form is special to dimension 2.
Now the first thing to do is to compute the partition function $`z(\lambda ,\mu )`$ or, equivalently, the free energy $`f(\lambda ,\mu )`$
$$z(\lambda ,\mu )=_{\text{states}}e^E,f(\lambda ,\mu )=\mathrm{ln}z(\lambda ,\mu )=\underset{g0}{}_{\text{metrics}}e^E.$$
The free energy is the sum of contributions of connected surfaces, while the partition function is the sum of contributions of all surfaces.
Neither of the above integrals is well-defined mathematically, but we would still like to compute them. To do that, physicists introduced a discrete model of Riemannian metrics, replacing them by quadrangulations (see also , Chapter 3 for a mathematical description). In this model, instead of considering Riemannian metrics, one considers metrics obtained by gluings of squares of area $`\epsilon `$. Our goal is to show that the ramified coverings of the sphere provide a new (maybe more natural) discrete model of Riemannian metrics.
Fix a positive number $`\epsilon `$. Consider a sphere with the standard (round) Riemannian metric of total area $`\epsilon `$. On this sphere, choose at random $`2n+2g2`$ points. Now chose a random connected $`n`$-sheeted covering of the sphere with simple ramifications over the $`2n+2g2`$ chosen points. The covering surface $`S`$ will automatically be of genus $`g`$. The metric on the sphere can be lifted to $`S`$, which will give us a metric with constant positive curvature except at the critical points, where it has conical singularities with angles $`4\pi `$. This metric is, of course, not Riemannian. However, one can argue that if $`\epsilon `$ is very small and the number of sheets very large, a random metric obtained in this way looks similar to a random Riemannian metric (unless we look at them through a microscope to reveal the difference). We do not know any rigorous statement that would formalize this intuitive explanation, but the same argument is used by physicists to justify the usage of quadrangulations.
Using our discrete model of metrics, one can write the free energy for the 2-dimensional gravity in the following way:
$$f(\lambda ,\mu )=\underset{g,n}{}\frac{\epsilon ^{2n+2g2}}{(2n+2g2)!}h_{g,n;\mathrm{}}e^{\lambda n\epsilon \mu (22g)}.$$
Here the factor $`\epsilon ^{2n+2g2}/(2n+2g2)!`$ is the volume of the space of choices of $`2n+2g2`$ unordered points on the sphere of area $`\epsilon `$, while $`n\epsilon `$ in the exponent is the area of the covering surface.
In Section 4.2 we show that if we let $`n\mathrm{}`$ while $`g`$ remains fixed, we have
$$\frac{h_{g,n;\mathrm{}}}{(2n+2g2)!}e^nn^{\frac{5}{2}(g1)1}b_g$$
for some constants $`b_g`$. Thus the coefficients of $`f`$ have the following asymptotic:
$$\frac{\epsilon ^{2n+2g2}}{(2n+2g2)!}h_{g,n;\mathrm{}}e^{\lambda n\epsilon \mu (22g)}b_ge^{(\lambda \epsilon 2\mathrm{ln}\epsilon +1)n}(\epsilon e^\mu )^{2g2}n^{\frac{5}{2}(g1)1}.$$
Now we make the final step by letting $`\epsilon `$ tend to $`0`$ in the expression of $`f`$. To obtain an interesting limit for the free energy, we must make $`\lambda `$ and $`\mu `$ depend on $`\epsilon `$. We want to use
$$\underset{n1}{}n^{\gamma 1}e^{\delta n}\frac{\mathrm{\Gamma }(\gamma )}{\delta ^\gamma }\text{ as }\delta 0.$$
Therefore we set $`\gamma =\frac{5}{2}(g1)`$ and we let $`\delta =\lambda \epsilon 2\mathrm{ln}\epsilon 1`$ tend to $`0`$, while
$$y=\frac{(\epsilon e^\mu )^{4/5}}{\delta }=\frac{(\epsilon e^\mu )^{4/5}}{\lambda \epsilon 2\mathrm{ln}\epsilon 1}$$
remains fixed. This gives us the final expression of the free energy, now depending on only one variable $`y`$:
$$f(y)=\mathrm{\Gamma }(5/2)b_0y^{5/2}b_1\mathrm{ln}y+\underset{g2}{}\mathrm{\Gamma }\left(5(g1)/2\right)b_gy^{5(1g)/2}.$$
The coefficients $`\mathrm{\Gamma }(5(g1)/2)b_g`$ are rational for odd $`g`$ and rational multiples of $`\sqrt{2}`$ for even $`g`$.
Our above treatment is parallel to E. Wittenโs treatment of the quadrangulation model in . Denote by $`Q_{g,n}`$ the number of ways to divide a surface of genus $`g`$ into $`n`$ squares. Then the study of the quadrangulation model involves the asymptotic of $`Q_{g,n}`$, which is given by
$$Q_{g,n}12^nn^{\frac{5}{2}(g1)1}b_g^{},$$
for another sequence of constants $`b_g^{}`$. This sequence was studied using matrix integrals, and it is known that a generating function for the sequence $`b_g^{}`$ satisfies the Painlevรฉ I equation. In the next section we show a similar result for the constants $`b_g`$. This implies that the functions $`f`$ obtained in the two models coincide up to a rescaling of the variable $`y`$; more precisely, we have
$$b_g^{}=2^{\frac{3}{2}(g1)+1}b_g.$$
In the treatment of the quadrangulation model in , Witten also introduced observables that correspond to counting quadrangulations with โimpuritiesโ, that is, the number of ways to divide a surface into a large number of squares and a fixed number of given polygons. Each observable $`\tau _d`$ is represented by a formal linear combination of a $`2`$-gon, a $`4`$-gon, and so on, up to a $`(2d+2)`$-gon. The values of these observables combine into a generating function $`F(t_0,t_1,\mathrm{})`$ that can be studied using matrix integrals. It turns out that $`^2F/t_0^2`$ is a solution of the Kortewegโde Vries (KdV) hierarchy. This solution is called the โstring solutionโ.
From now on the notation $``$ or $`_g`$ will mean
$$\tau _{d_1}\mathrm{}\tau _{d_n}_g=\tau _{d_1}\mathrm{}\tau _{d_n}=\underset{\overline{}_{g,n}}{}\psi _1^{d_1}\mathrm{}\psi _n^{d_n}.$$
It turns out that the generating series
$$F(t_0,t_1,\mathrm{})=\underset{n1}{}\underset{d_1,\mathrm{},d_n}{}\tau _{d_1}\mathrm{}\tau _{d_n}\frac{t_{d_1}\mathrm{}t_{d_n}}{n!}.$$
coincides with the series $`F`$ obtained from the quadrangulation model. In particular, its second derivative $`U=^2F/t_0^2`$. satisfies the KdV equation:
$$\frac{U}{t_1}=U\frac{U}{t_0}+\frac{1}{12}\frac{^3U}{t_0^3}.$$
(1)
This was conjectured by E. Witten in and proved by M. Kontsevich in .
In Section 4.3 we show that the numbers $`\tau _{d_1}\mathrm{}\tau _{d_n}`$ can also be obtained in the model of ramified coverings. Each observable $`\tau _d`$ is represented by a formal linear combination of a noncritical point, a simple critical point, and so on, up to a $`d`$-tuple critical point.
Recently, M. Kazarian and S. K. Lando found an independent proof of the fact that the function $`U`$ arising in the enumeration of Hurwitz numbers satisfies the KdV equation. This has lead them to a new proof of Wittenโs conjecture.
### 4.2 The Painlevรฉ I equation
The results of this section were obtained in common with S. Lando.
###### Proposition 4.1
For a fixed $`g`$, we have
$$\frac{h_{g,n;\mathrm{}}}{(2n+2g2)!}e^nn^{\frac{5}{2}(g1)1}b_g\text{as}n\mathrm{},$$
where
$$b_g=\frac{\tau _2^{3g3}}{(3g3)!\mathrm{\hspace{0.33em}\hspace{0.33em}2}^{\frac{5}{2}(g1)}\mathrm{\Gamma }\left(\frac{5}{2}(g1)\right)}.$$
#### Proof.
From Theorem 5 we see that $`H_{g;\mathrm{}}`$ is a polynomial in $`Z`$ with leading term
$$\frac{\tau _2^{3g3}}{(3g3)!}Z^{5g5}.$$
Indeed, if $`l=3g3`$, the degree of the class $`\psi _1^2\mathrm{}\psi _l^2`$ is exactly the dimension of $`\overline{}_{g,l}`$ (both are equal to $`6g6`$). Therefore the classes $`c(W^{})`$ and $`1/(1\psi _i)`$ do not contribute. On the other hand, by Proposition 2.7 the coefficients of the series $`Z^l`$ grow as
$$\frac{n^{\frac{l}{2}1}e^n}{\mathrm{\Gamma }(l/2)\mathrm{\hspace{0.33em}2}^{l/2}}.$$
Multiplying the coefficient of the leading term
$$\frac{\tau _2^{3g3}}{(3g3)!}Z^{5g5}$$
of $`H_g`$ by the asymptotic of coefficients of $`Z^{5g5}`$ we obtain the leading term of the asymptotic of $`h_{g,n}`$, in particular, the constant $`b_g`$. โ
The first values of the constants $`b_g`$ are
$$b_0=\frac{1}{\sqrt{2\pi }},b_1=\frac{1}{2^43},b_2=\frac{1}{\sqrt{2\pi }}\frac{7}{2^53^35},$$
$$b_3=\frac{57^2}{2^{16}3^5},b_4=\frac{1}{\sqrt{2\pi }}\frac{75297}{2^{11}3^85^21113}.$$
The above expression for $`b_g`$ allows us to rewrite the function $`f^{\prime \prime }(y)=u(y)`$ in the following way:
$$u(y)=\sqrt{2y}+\frac{1}{12}(2y)^2+\underset{g2}{}(55g)(35g)\frac{\tau _2^{3g3}}{(3g3)!}(2y)^{\frac{15g}{2}}.$$
(2)
###### Proposition 4.2
The function $`u(y)`$ satisfies the Painlevรฉ I equation
$$u^2(y)+\frac{1}{6}u^{\prime \prime }(y)=2y.$$
#### Proof.
By extracting the coefficient of $`t_2^{3g1}`$ in the KdV equation (1), we obtain, for every $`g1`$,
$$\frac{\tau _0^2\tau _1\tau _2^{3g1}_g}{(3g1)!}=\underset{\stackrel{g^{}+g^{\prime \prime }=g}{g^{}1,g^{\prime \prime }0}}{}\frac{\tau _0^2\tau _2^{3g^{}1}_g^{}}{(3g^{}1)!}\frac{\tau _0^3\tau _2^{3g^{\prime \prime }}_{g^{\prime \prime }}}{(3g^{\prime \prime })!}+\frac{1}{12}\frac{\tau _0^5\tau _2^{3g1}_{g1}}{(3g1)!}.$$
(3)
Now we use the string and the dilaton equations to kill the $`\tau _0`$ and $`\tau _1`$ factors in all the brackets of (3). We obtain the following identities (with some exceptions for low genus):
$`{\displaystyle \frac{\tau _0^2\tau _1\tau _2^{3g1}_g}{(3g1)!}}`$ $`=`$ $`(5g5)(5g3)(5g1){\displaystyle \frac{\tau _2^{3g3}_g}{(3g3)!}},{\displaystyle \frac{\tau _0^2\tau _1\tau _2^2_1}{2!}}={\displaystyle \frac{1}{3}}.`$
$`{\displaystyle \frac{\tau _0^2\tau _2^{3g^{}1}_g^{}}{(3g^{}1)!}}`$ $`=`$ $`(5g5)(5g3){\displaystyle \frac{\tau _2^{3g^{}3}_g^{}}{(3g^{}3)!}},{\displaystyle \frac{\tau _0^2\tau _2^2_1}{2!}}={\displaystyle \frac{1}{12}}.`$
$`{\displaystyle \frac{\tau _0^3\tau _2^{3g^{\prime \prime }}_{g^{\prime \prime }}}{(3g^{\prime \prime })!}}`$ $`=`$ $`(5g5)(5g3)(5g1){\displaystyle \frac{\tau _2^{3g^{\prime \prime }3}_{g^{\prime \prime }}}{(3g^{\prime \prime }3)!}},`$
$`\tau _0^3_0=1,{\displaystyle \frac{\tau _0^3\tau _2^3_1}{3!}}={\displaystyle \frac{1}{3}}.`$
$`{\displaystyle \frac{\tau _0^5\tau _2^{3g1}_{g1}}{(3g1)!}}`$ $`=`$ $`(5g10)(5g8)(5g6)(5g4)(5g2){\displaystyle \frac{\tau _2^{3g6}_{g1}}{(3g6)!}},`$
$`{\displaystyle \frac{\tau _0^5\tau _2^2_0}{2!}}=3,{\displaystyle \frac{\tau _0^5\tau _2^5_1}{5!}}=16.`$
Substitute these expressions in (3) and compare to the expression (2) of $`u`$. Taking into account the exceptional starting terms we obtain
$$\sqrt{2y}\left(u^{}(y)+\frac{1}{\sqrt{2y}}\right)=\left(u(y)+\sqrt{2y}\right)u^{}(y)+\frac{1}{12}u^{\prime \prime \prime }(y).$$
We rewrite this as
$$u(y)u^{}(y)+\frac{1}{12}u^{\prime \prime \prime }(y)=1$$
and integrate it once to obtain
$$u^2(y)+\frac{1}{6}u^{\prime \prime }(y)=2y.$$
### 4.3 The KdV hierarchy
Now we show how to obtain the numbers $`\tau _{d_1}\mathrm{}\tau _{d_p}`$ as leading term coefficients of the asymptotic of Hurwitz numbers.
Note that Okounkov and Pandharipande also obtained the numbers $`\tau _{d_1}\mathrm{}\tau _{d_p}`$ using asymptotics of Hurwitz numbers. However their asymptotics are different from ours and have no direct physical interpretation.
The Hurwitz numbers involved are those that count ramified coverings with many simple ramification points, but only one multiple ramification point that has $`p`$ marked preimages. Each marked preimage corresponds to a factor $`\tau `$ in the bracket. Using the notation from Section 3, the numbers we are interseted in are $`|\mathrm{Aut}(\mu )|h_{g,n;\mu }`$, with $`\mu =(b_1,\mathrm{},b_p)`$. The factor $`|\mathrm{Aut}(\mu )|`$ is due to the fact that the marked preimages are numbered.
Denote by b a $`b`$-tuple preimage of the special ramification point. Our result is then best described by the following symblic formula:
$$\tau _d\text{}=\text{}\frac{1}{0!(d+1)^d}\frac{\overline{)d+1}}{Y^{d+1}}\frac{1}{1!d^{d1}}\frac{d}{Y^d}+\mathrm{}+(1)^d\frac{1}{d!\mathrm{\hspace{0.17em}1}^0}\frac{1}{Y}$$
The recipe for obtaining the number $`\tau _{d_1}\mathrm{}\tau _{d_p}`$ is the following.
1. Replace each $`\tau _d`$ by the right-hand side of the above symblic equality.
2. Expand the product to obtain a linear combination of terms of the form
$$\text{const}\frac{\overline{)b_1}\mathrm{}\overline{)b_p}}{Y^{b_1+\mathrm{}+b_p}}.$$
3. To each such term assign the partition $`\mu =(b_1,\mathrm{},b_p)`$ and the series
$$\text{const}\frac{|\mathrm{Aut}(\mu )|H_{g;\mu }}{Y^{b_1+\mathrm{}+b_p}}.$$
4. Add all the series thus obtained. This gives a series in $`๐`$ that we denote by $`H[\tau _{d_1}\mathrm{}\tau _{d_p}](q).`$
###### Theorem 6
We have
$$H[\tau _{d_1}\mathrm{}\tau _{d_p}]=\tau _{d_1}\mathrm{}\tau _{d_p}(Z+1)^{2g2+l}.$$
The asymptotic of the coefficient of $`q^n`$ (as $`n\mathrm{}`$) in $`H[\tau _{d_1}\mathrm{}\tau _{d_p}](q)`$ equals
$$\frac{\tau _{d_1}\mathrm{}\tau _{d_p}}{2^{\frac{2g2+p}{2}}\mathrm{\Gamma }\left(\frac{2g2+p}{2}\right)}e^nn^{\frac{2g2+p}{2}1}.$$
#### Proof.
Here again we will use Theorem 5. The crucial part of this proposition is the polynomial
$$\phi _{b_1,\mathrm{},b_p}(Z)=\underset{l0}{}\frac{Z^l}{l!}\underset{\overline{}_{g,p+l}}{}\frac{1}{(1b_1\psi _1)\mathrm{}(1b_p\psi _p)}\frac{c(W^{})\psi _{p+1}^2\mathrm{}\psi _{p+l}^2}{(1\psi _{p+1})\mathrm{}(1\psi _{p+l})}.$$
We are going to consider linear combinations of such polynomials for different $`b_i`$โs. Our goal is to obtain a cancellation of all higher order terms in $`Z`$ leaving only a constant term ($`l=0`$). This constant term will turn out to be $`\tau _{d_1}\mathrm{}\tau _{d_p}`$.
First of all, here is a linear combination of the series $`1/(1\psi )`$, $`1/(12\psi )`$, โฆ, $`1/(1(d+1)\psi )`$ whose terms up to $`\psi ^{d1}`$ vanish:
$$\frac{1}{d!}\underset{b=1}{\overset{d+1}{}}\frac{(1)^{d+1b}\left(\genfrac{}{}{0pt}{}{d}{b1}\right)}{1b\psi }=\psi ^d+O(\psi ^{d+1}).$$
(4)
In the expression
$$|\mathrm{Aut}(\mu )|H_{g;\mu }=(Z+1)^{2g2+p}Y^{b_1+\mathrm{}+b_p}\underset{i=1}{\overset{p}{}}\frac{b_i^{b_i}}{b_i!}\phi (Z),$$
the integrand $`1/(1b_i\psi _i)`$ appears with an additional factor $`Y^{b_i}b_i^{b_i}/b_i!`$ in front of the integral. To compensate for this factor, we multiply the coefficients of (4) by its inverse, which gives
$$\underset{b=1}{\overset{d+1}{}}\frac{(1)^{d+1b}}{(d+1b)!b^{b1}}\frac{\text{b}}{Y^b}.$$
This is precisely the formula that we gave for $`\tau _d`$.
Multiplying such expressions for $`d=d_1,\mathrm{},d_p`$ and adding them up we obtain
$$H[\tau _{d_1}\mathrm{}\tau _{d_n}]=$$
$$(Z+1)^{2g2+p}\underset{l0}{}\frac{Z^l}{l!}\underset{\overline{}_{g,p+l}}{}\left(\psi _1^{d_1}\mathrm{}\psi _p^{d_p}+\text{h.o.t.}\right)\frac{c(W^{})\psi _{p+1}^2\mathrm{}\psi _{p+l}^2}{(1\psi _{p+1})\mathrm{}(1\psi _{p+l})},$$
where โh.o.t.โ means โhigher order termsโ.
The above integral vanishes for dimension reasons whenever $`l>0`$. For $`l=0`$, the factor $`c(W^{})`$ contributes only by $`c_0(W^{})=1`$. Thus $`H[\tau _{d_1}\mathrm{}\tau _{d_n}]=\tau _{d_1}\mathrm{}\tau _{d_p}(Z+1)^{2g2+p}`$ as claimed.
The second assertion of the theorem follows from the first one and from the asymptotic of the coefficients of $`Z^k`$ (Proposition 2.7). โ
###### Remark 4.3 (Kazarian)
Equality (4) can be used to express the numbers $`\tau _{d_1}\mathrm{}\tau _{d_n}`$ as finite linear combinations of Hurwitz numbers, without considering any asymptotics. More precisely, we have
$$\tau _{d_1}\mathrm{}\tau _{d_n}=\underset{b_1,\mathrm{},b_p}{}\left(\underset{i=1}{\overset{p}{}}\frac{(1)^{d_i+1b_i}}{(d_i+1b_i)!b_i^{b_i1}}\right)\frac{|\mathrm{Aut}(b_1,\mathrm{},b_p)|h_{g,n;b_1,\mathrm{},b_p}}{c(n)!}.$$
Here the sum is over $`1b_id_i+1`$, the number of sheets is $`n=b_i`$ and $`c(n)=n+p+2g2`$.
### 4.4 Ramified coverings of a torus and abelian differentials
Fix an integer $`g1`$ and a list of $`p`$ nonnegative integers $`b_1,\mathrm{},b_p`$ with the condition $`b_i=2g2`$. We consider the space $`D_{g;b_1,\mathrm{},b_p}`$ of abelian differentials on Riemann surfaces of genus $`g`$, with zeroes of multiplicities $`b_1,\mathrm{},b_p`$. More precisely, $`D_{g;b_1,\mathrm{},b_p}`$ is the space of triples $`(C,\{x_1,\mathrm{},x_p\},\alpha )`$, where $`C`$ is a smooth complex curve, $`x_1,\mathrm{},x_pC`$ are distinct marked points, and $`\alpha `$ is an abelian ($`=`$ holomorphic) differential on $`C`$ whose zero divisor is precisely $`b_1x_1+\mathrm{}+b_px_p`$.
It turns out that the space $`D_{g;b_1,\mathrm{},b_p}`$ has a natural integer affine structure. This means that it can be covered by charts of local coordinates in such a way that the transition functions are affine maps with integer coefficients. Such local coordinates are introduced as follows. Fix a basis $`l_1,\mathrm{},l_{2g+p1}`$ of the relative homology group $`H_1(C,\{x_1,\mathrm{},x_p\},)`$. Then the integrals of $`\alpha `$ over the cycles $`l_i`$ are the local coordinates we need. The area function
$$A:(C,\alpha )\frac{i}{2}_C\alpha \overline{\alpha }$$
is a quadratic form with respect to the affine structure.
The integer affine structure allows one to define a volume measure on the space $`D_{g;b_1,\mathrm{},b_p}`$. It is then a natural question to find the total volume of the part of the space $`D_{g;b_1,\mathrm{},b_p}`$ defined by $`A1`$ (the volume of the whole space being infinite).
A. Eskin and A. Okounkov obtained an effective way to calculate these volumes using the asymptotic for the number of ramified coverings of a torus. Consider the elliptic curve obtained by gluing the opposite sides of the square $`(0,1,i,1+i)`$ endowed with the abelian differential $`dz`$. Given a ramified covering of this elliptic curve with critical points of multiplicities $`b_1,\mathrm{},b_p`$, we can lift the abelian differential to the covering curve and obtain a point of $`D_{g;b_1,\mathrm{},b_p}`$. One can then easily show that such points are densely and uniformly distributed in $`D_{g;b_1,\mathrm{},b_p}`$ if one considers coverings with a big number of sheets. Moreover, R. Dijkgraaf and S. Bloch and A. Okounkov showed that the generating series for the ramified coverings of the torus that arise in this study are quasi-modular forms. In other words, they lie in the algebra
$$[E_2,E_4,E_6],$$
where $`E_{2k}`$ are the Eisenstein series
$$E_{2k}(q)=\frac{1}{2}\zeta (12k)+\underset{n1}{}\left(\underset{d|n}{}d^{2k1}\right)q^n.$$
We conclude with the following comparison between the counting of ramified coverings of a sphere and of a torus.
Sphere: The generating series enumerating the ramified coverings lie in the algebra $`๐`$.
Torus: The generating series enumerating the ramified coverings lie in the algebra of quasi-modular forms.
Sphere: The coefficients of a generating series grow as $`e^nn^{\gamma 1}c`$. The exponent $`\gamma `$ is a half-integer. The number $`\gamma `$ is called the string susceptiblity. The constant $`c`$ is an observable in 2-dimensional gravity.
Torus: The sum of the first $`n`$ coefficients of a generating series grows as $`n^dc`$. The number $`d`$ is the complex dimension of the corresponding space of abelian differentials. The constant $`c`$ is its volume.
Sphere: The observables can be arranged into a generating series that is a solution of the KdV hierarchy.
Torus: As far as we know, nobody has tried to arrange the volumes of the spaces of abelian differentials into a unique generating series.
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# Diabatic limit, eta invariants and Cauchy-Riemann manifolds of dimension 3
## 1. Introduction
In the first two authors of this paper introduced a new invariant, called the $`\nu `$-invariant, of strictly pseudoconvex Cauchy-Riemann (CR) compact 3-manifolds. This invariant was obtained by taking the limit of the $`\eta `$-invariants of an adequately defined (but quite complicated) sequence of Riemannian metrics approximating the CR structure, after cancellation of the possibly diverging terms by adding well-chosen local contributions. We claimed in that this invariant was an analogue in CR geometry of the $`\eta `$-invariant in conformal geometry. However, its rather abstract definition makes it difficult to compute explicit expressions for it or to get a further understanding of its properties. The goal of this paper is then to provide links between $`\nu `$ and other natural $`\eta `$-invariants in CR geometry.
In a first step, we introduce a renormalized $`\eta `$-invariant that takes into account the fact that CR geometry can be seen as a limit of a sequence of conformal structures that diverges outside the contact distribution. If a compatible contact form $`\theta `$ is fixed on the CR manifold $`M`$, one considers the family of metrics
(1)
$$h_\epsilon =\epsilon ^1\theta ^2+\gamma ,$$
where $`\gamma =d\theta (,J)`$ and $`J`$ is the underlying complex structure on the contact distribution. *When $`\epsilon `$ goes to $`0`$* the metrics $`h_\epsilon `$ blow up except in the contact distribution, and therefore diverge to the Carnot-Carathรฉodory metric associated to the CR structure and the contact form (this is one of the main motivation for considering this kind of sequences). A natural object one can consider is the constant term $`\eta _0`$ in an asymptotic expansion for $`(\eta (h_\epsilon ))`$ in powers of $`\epsilon `$, *when $`\epsilon `$ goes to $`0`$*. This always exists, as we shall see, and we shall call it the *renormalized $`\eta `$-invariant of the pseudohermitian manifold $`(M,\theta )`$*. This invariant is of course much more easily studied than the $`\nu `$-invariant, because it is built from the sequence (1) of metrics that is much simpler than the one used to build $`\nu `$ in . Note however that it is a pseudohermitian invariant, *i.e.* it depends on the choice of $`\theta `$, contrarily to $`\nu `$.
In the other direction, *i.e. when $`\epsilon `$ goes to $`\mathrm{}`$*, on can also obtain another natural invariant in case the Tanaka-Webster torsion of $`(M,\theta )`$ vanishes, that is when the action of the Reeb vector field is isometric. In this case, $`\eta (h_\epsilon )`$ converges and its limit $`\eta _{\mathrm{ad}}`$ is the so called *adiabatic limit*. It has attracted much attention in the past few years, see for instance. We shall call the reverse process of taking a limit when $`\epsilon `$ goes to $`0`$ a *diabatic* limit. When torsion vanishes, it turns out that the diabatic $`\eta _0`$ equals the adiabatic $`\eta _{\mathrm{ad}}`$.
Our first result shows that the difference between the CR invariant $`\nu `$ and the pseudo-hermitian $`\eta _0`$ is an integral of a local contribution involving the square of the Tanaka-Webster curvature.
###### 1.1 Theorem.
For any compact strictly pseudoconvex Cauchy-Riemann $`3`$-manifold $`M`$, and any choice $`\theta `$ of contact form, one has
(2)
$$\nu (M)=\mathrm{\hspace{0.17em}3}\eta _0(M,\theta )+\frac{1}{16\pi ^2}_MR^2\theta d\theta ,$$
where $`R`$ is the Tanaka-Webster curvature of $`(M,\theta )`$.
This yields a new definition of the $`\nu `$-invariant, see Remark 4.2, together with some explicit computations: they can be done on manifolds on which $`\eta _0`$ is computable. We are then able to apply this to transverse $`๐^1`$-invariant CR structures on Seifert manifolds. The CR manifolds we are interested in come with a locally free action of $`๐^1`$ that is transverse to the contact distribution, and preserves both the contact and the complex structures. We shall call them *Cauchy-Riemann-Seifert manifolds* (in short CR-Seifert). We refer to for more information on the more general class of $`๐^1`$-invariant CR structures. CR-Seifert manifolds can be efficiently described as orbifold $`๐^1`$-bundles over $`2`$-dimensional orbifolds. At each orbifold point on the base, the orbifold bundle data consists of the following: if the local fundamental group is $`/\alpha `$ ($`\alpha ^{}`$), a generator acts on a local chart around $`p`$ on the basis manifold as $`\mathrm{e}^{i\frac{2\pi }{\alpha }}`$ and on the fiber as $`\mathrm{e}^{i\frac{2\pi \beta }{\alpha }}`$ with $`\beta `$ prime to $`\alpha `$. The orbifold $`๐^1`$-bundles are topologically classified by their degrees (first Chern numbers), which are in this case rational numbers. One then endows the manifold with an invariant strictly pseudoconvex CR structure as follows: the underlying contact structure is provided by an equivariant connection $`1`$-form on the bundle, whereas the complex structure is induced from the basis (orbifold) Riemann surface; the strict pseudoconvexity condition constrains the degrees $`d`$ of these $`๐^1`$-bundles to be negative.
###### 1.2 Theorem.
Let $`M`$ be a compact strictly pseudoconvex CR-Seifert $`3`$-manifold, of degree $`d`$ over the orbifold surface $`\mathrm{\Sigma }`$, and with $`๐^1`$-action generated by the Reeb field of a contact form $`\theta `$. If $`R`$ is the Tanaka-Webster curvature of $`(M,\theta )`$, then
(3)
$$\nu (M)=d312\underset{j=1}{\overset{p}{}}s(\alpha _j,1,\beta _j)+\frac{1}{8\pi }_\mathrm{\Sigma }R^2๐\theta ,$$
where $`s(\alpha ,\rho ,\beta )`$ is the Rademacher-Dedekind sum $`\frac{1}{4\alpha }\underset{k=1}{\overset{\alpha 1}{}}\mathrm{cot}\left(\frac{k\rho \pi }{\alpha }\right)\mathrm{cot}\left(\frac{k\beta \pi }{\alpha }\right)`$ .
The Tanaka-Webster curvature $`R`$ of such an $`(M,\theta )`$ actually coincides with Riemannian curvature of the base $`\mathrm{\Sigma }`$, if it is endowed with the metric $`\gamma =d\theta (,J)`$. When this curvature is constant, (3) specializes into the following interesting formula, which shows that the $`\nu `$-invariant is a topological invariant in this case:
###### 1.3 Corollary.
Let $`M`$ be a CR-Seifert manifold as above, with constant Tanaka-Webster curvature. Let $`\chi `$ be the rational Euler characteristic of $`\mathrm{\Sigma }`$. Then,
(4)
$$\nu (M)=d3\frac{\chi ^2}{4d}12\underset{j=1}{\overset{p}{}}s(\alpha _j,1,\beta _j).$$
However, Theorem 1.1 is not entirely satisfactory, as it provides a link between the CR invariant $`\nu `$ and the diabatic invariant $`\eta _0`$; one would instead prefer a relationship between $`\nu `$ and invariants defined directly in terms of the CR or pseudohermitian geometry. One such object is the contact-de Rham complex introduced in , and especially the $`\eta `$-invariant of the middle degree operator appearing there.
The relevant operator (denoted by $`D`$ henceforth) is the analogue in this setting of the boundary operator for the signature $`\pm (dd)`$ that gives rise to the $`\eta `$-invariant on $`3`$-dimensional Riemannian manifolds. It is known that the spectrum of the operator $`D`$ appears in the rescaled limit of the collapsing spectrum of $`P_\epsilon =\pm (d_\epsilon _\epsilon d)`$ for the metrics $`h_\epsilon `$ of (1), when performing the diabatic limit . However, this limit is not uniform enough to yield a direct relation between the $`\eta `$-invariants. In this paper, we prove a general relation between $`\nu `$ and $`\eta (D)`$ in the special case provided by our transverse $`๐^1`$-invariant CR manifolds. In effect, we show that $`\eta (D)`$ and $`\nu `$ differ only by a simple local term in the Tanaka-Webster curvature of any chosen pseudohermitian structure. Our second main set of results then reads:
###### 1.4 Theorem.
Let $`M`$ be a compact strictly pseudoconvex CR-Seifert $`3`$-manifold, with $`๐^1`$-action generated by the Reeb field of a contact form $`\theta `$. If $`R`$ is the Tanaka-Webster curvature of $`(M,\theta )`$ and $`D`$ is the middle operator of the contact complex, then
(5)
$$\eta _0(M,\theta )=\eta (D)+\frac{1}{512}_MR^2\theta d\theta .$$
###### 1.5 Corollary.
Let $`M`$ be a CR-Seifert $`3`$-manifold as above, then one has:
(6)
$$\nu (M)=\mathrm{\hspace{0.17em}3}\eta (D)+(\frac{1}{16\pi ^2}\frac{3}{512})_MR^2\theta d\theta .$$
The philosophy underlying our results is indeed the following: whereas $`\nu `$ is easily related to $`\eta _0`$, $`\eta (D)`$ compares itself more easily with $`\eta _0`$ rather than to $`\nu `$. This somehow โexplainsโ the quite strange combination of constants appearing in front of the curvature term in (6) in theorem 1.5: it is a sum of diabatic contribution steming from theorem 1.1 and a purely spectral term linking $`\eta (D)`$ and $`\eta _0`$, as will be apparent from section 7.
For general CR manifolds, we expect that when we take the diabatic limit $`\epsilon 0`$, the collapsing spectrum of $`P_\epsilon `$ gives the contribution $`\eta (D)`$ in the limit, while the remaining part of the spectrum, after renormalization, gives only an integral of local terms. This leads to the following conjecture.
###### 1.6 Conjecture.
There exists a constant $`C`$ such that, for any compact strictly pseudoconvex Cauchy-Riemann $`3`$-manifold $`M`$ and any choice $`\theta `$ of contact form, one has
(7)
$$\nu (M)=\mathrm{\hspace{0.17em}3}\eta (D)+(\frac{1}{16\pi ^2}\frac{3}{512})_MR^2\theta d\theta +C_M|\tau |^2\theta d\theta ,$$
with $`R`$ and $`\tau `$ the Tanaka-Webster curvature and torsion of $`(M,\theta )`$.
As a first indication for the conjecture, we shall give in Theorem 9.4 an abstract argument that shows that there exists a CR invariant of the form $`\eta (D)+C_1R^2+C_2|\tau |^2`$. Unfortunately, we are unable to calculate the constants completely, see Remark 9.6.
It is known that the $`\eta `$-invariant of the boundary operator for signature is conformally invariant. If the conjecture is true, then this is no more the case for $`\eta (D)`$, which is a priori an invariant of the pseudo-hermitian structure only: it depends on the choice of a metric in the conformal class adapted to the CR structure.
A third goal of this paper is to provide some geometric applications on CR-Seifert manifolds, mainly with constant curvature. They are spherical (locally isomorphic to the standard CR sphere $`๐^3`$), hence are the boundary at infinity of a complex hyperbolic metric defined in a neighbourhood $`(0,\epsilon )\times M`$ of $`M`$ (in the case of the $`3`$-sphere we can of course extend the metric globally to get the Bergmann metric on the 4-ball). From \[11, Theorem 1.2\] and Theorem 1.3, we get the following obstruction for this neighbourhood to have a global extension to a smooth complex hyperbolic surface (with only one end):
###### 1.7 Corollary.
If a CR-Seifert manifold $`M^3`$ is the boundary at infinity of a complex hyperbolic metric defined on the interior of a smooth compact manifold $`N^4`$ with boundary $`M`$, then one has necessarily $`\nu (M)=\chi (N)+3\tau (N)`$, where $`\chi (N)`$ and $`\tau (N)`$ denote the Euler characteristic and signature of $`N`$. In particular, $`\nu (M)`$, as provided by the formula *(3)*, is an integer.
This is a topological constraint on a filling, which we can restate in the smooth case (no orbifold singularities):
###### 1.8 Corollary.
Let $`M`$ be a $`๐^1`$-bundle of degree $`d`$ over a Riemann surface $`\mathrm{\Sigma }`$ of Euler characteristic $`\chi `$, with a $`๐^1`$-invariant spherical CR structure. If $`\frac{\chi ^2}{4d}`$ is not an integer then $`M`$ is not the boundary at infinity of a complex hyperbolic metric.
The case $`d=\frac{\chi }{2}`$ yields an integer, and indeed, if $`\mathrm{\Sigma }`$ is hyperbolic, $`N`$ can be taken to be the disk bundle of a square root of the tangent bundle of $`\mathrm{\Sigma }`$, which is well known to carry a complex hyperbolic metric issued from a representation of $`\pi _1(\mathrm{\Sigma })`$ in $`SU(1,1)SU(1,2)`$. Our obstruction then gives an interesting hint on whether a spherical CR-Seifert $`3`$-manifold may appear as the quotient of the complement of the limit set in the $`3`$-sphere of some discrete fixed point-free subgroup of $`SU(1,2)`$ .
More generally, the calculation in Theorem 1.2 gives an obstruction for $`M`$ to be the boundary at infinity of a Kรคhler-Einstein or Einstein metric. The manifolds considered in this paper are known to bound a complex Stein space with at most a finite number of singular points and one may wish to endow it with a Kรคhler-Einstein metric as in Cheng-Yau . The type of metric to be considered has the same kind of asymptotic expansion near the boundary $`M`$ as the Bergman metric ; we called them โasymptotically complex hyperbolicโ (ACH) in . If no singular points are present and if the Cheng-Yau metric exists, one gets from the Miyaoka-Yau inequality proved in the following:
###### 1.9 Corollary.
Let $`M`$ be as in Theorem 1.2. If $`M`$ is the boundary at infinity of an ACH Einstein metric on $`M^4`$, such that a Kronheimer-Mrowka invariant of $`(N,M)`$ is nonzero (in particular, if $`M`$ is the boundary at infinity of a Kรคhler-Einstein metric on $`N`$), then
$$\chi (N)3\tau (N)\nu (M)=d+3+12\underset{j=1}{\overset{p}{}}s(\alpha _j,1,\beta _j)\frac{1}{8\pi }_\mathrm{\Sigma }R^2๐\theta .$$
For more information on Stein fillings, see . The Kronheimer-Mrowka invariants are Seiberg-Witten type invariants defined for a compact 4-manifold with contact boundary; in particular, they do not vanish if $`M`$ carries a symplectic form compatible with the contact structure on the boundary, and this implies the Miyaoka-Yau inequality . This inequality can of course be obtained directly for Kรคhler-Einstein metrics.
The paper is organized as follows. After recalling the definition of the $`\nu `$-invariant in section 2, we define the renormalized $`\eta `$-invariant $`\eta _0`$ and compare it to $`\nu `$ in sections 3 and 4. The proof relies on relatively simple considerations on $`\eta `$-invariants and Chern-Simons theory, that prove that the difference between $`\nu +3\eta _0`$ is necessarily of the expected form: an integral term in the square of the curvature and the squared norm of the torsion. The constants in front of these local terms are then computed by considering sufficiently many examples: left invariant structures on the $`3`$-sphere.
The reader will then find in section 5 the explicit computations of $`\nu `$ on CR-Seifert manifolds.
Taking one step further, sections 6 to 8 lead to the relation between $`\eta _0`$ and $`\eta (D)`$ in the case of transverse $`๐^1`$-invariant CR structures. The proof of Theorem 1.5 relies on an explicit study of the spectra of the $`D`$ operator and the boundary operator for the signature $`\pm (d_\epsilon d_\epsilon )`$ on closed $`2`$-forms for the sequence of Riemannian metrics $`h_\epsilon `$ that performs the diabatic limit in (1). This can be done only for $`๐^1`$-invariant structures and index theory shows once again that a relation of the expected type must exist. One then has again to evaluate the constant in front of the integral term by looking at explicit computations of both $`\eta (D)`$ and $`\nu `$ on the standard sphere.
The existence of a CR invariant of the form $`\eta (D)+C_1R^2+C_2|\tau |^2`$ is considered in section 9. We also present a proof of the existence of $`\eta (D)`$ on any compact strictly pseudoconvex CR manifold of dimension $`3`$, a fact certainly known to specialists but whose proof seems to have never been published so far.
The paper ends with a short section 10 devoted to the proof of the corollaries and to some generalizations, and also to a comparison with the results one can get in the Kรคhler-Einstein case using the $`\mu `$-invariant of Burns and Epstein .
## 2. The $`\nu `$-invariant
Let $`M`$ be a $`3`$-dimensional compact strictly pseudoconvex CR manifold, *i.e.* a compact manifold $`M`$ endowed with a complex structure $`J`$ defined on a contact distribution $`H`$ in $`TM`$.
A pseudohermitian structure $`(M,\theta )`$ consists in the additional choice of a contact form $`\theta `$. It induces a metric $`\gamma =d\theta (,J)`$ on $`H`$ and a splitting of both $`TM`$ and $`T^{}M`$ by means of the Reeb vector field $`T`$ defined by $`\theta (T)=1`$ and $`\iota _Td\theta =0`$. The Tanaka-Webster connection is then defined by working in a local coframe $`(\theta ,\theta ^1,\theta ^{\overline{1}})`$ such that $`d\theta =i\theta ^1\theta ^{\overline{1}}`$: the connection form is a purely imaginary 1-form $`\omega _1^1`$, and the torsion $`\tau ^1`$ is a $`(0,1)`$-form such that
$$d\theta ^1=\theta ^1\omega _1^1+\theta \tau ^1,$$
and the curvature $`R`$ is defined by
$$d\omega _1^1=iRd\theta +(\tau _{,\overline{1}}^{\overline{1}}\tau _{,1}^1)\theta .$$
In more invariant terms, it is the only metric and complex compatible connection $`^W`$ on $`H`$ such that the torsion $`\tau =T^^W(T,)_{|H}`$ anticommutes with $`J`$.
Given a pseudohermitian manifold $`(M,\theta )`$, one can define a first metric $`g_0`$ on the product $`N=_+\times M`$ by
(8)
$$g_0=dr^2+h_0(r),h_0(r)=\mathrm{e}^{2r}\theta ^2+\mathrm{e}^r\gamma .$$
Here we think of the initial $`M`$ as a boundary of $`M`$ at infinity (*i.e.* when $`r`$ goes to infinity). Remark that when doing a conformal change $`\theta ^{}=f\theta `$, one gets a metric $`g_0^{}=(dr^{})^2+\mathrm{e}^{2r^{}}f^2\theta ^2+\mathrm{e}^r^{}f\gamma `$, and the difference $`g_0^{}g_0`$ goes to zero at infinity after the coordinate change $`r=r^{}+\mathrm{log}f`$. Therefore the asymptotic behaviour of the metric $`g_0`$ depends only on the CR structure. We note moreover that
$$h_0(r)=\mathrm{e}^r(\mathrm{e}^r\theta ^2+\gamma )=\epsilon ^1h_\epsilon ,$$
where $`h_\epsilon `$ is the metric introduced in equation (1), with $`\epsilon =\mathrm{e}^r`$.
One can extend $`J`$, initially defined on $`M`$, to an almost complex structure $`J_0`$ on $`N`$ by defining
$$J_0_r=\mathrm{e}^rT,$$
where $`T`$ is the Reeb field associated to $`\theta `$. As explained in the curvature of $`g_0`$, together with $`J_0`$, is asymptotic when $`r+\mathrm{}`$ to curvature of the complex hyperbolic plane with holomorphic sectional curvature $`1`$.
One can add higher order corrections to $`J_0`$ and $`g_0`$ to get a uniquely defined jet of Kรคhler-Einstein metric $`g_{KE}`$ up to order $`\mathrm{e}^{2r}`$ (relatively to $`g_0`$), when $`r`$ tends to infinity. This development is expressed with the covariant derivatives of Tanaka-Webster curvature and torsion $`R`$ and $`\tau `$ of the pseudohermitian manifold $`(M,\theta )`$, and calculated in \[11, theorem 3.3 and corollary 3.4\]. More precisely, one finds an infinite series $`J(r)=J_0+J_1\mathrm{e}^r+J_2\mathrm{e}^{2r}+\mathrm{}`$ giving an integrable (formal) complex structure $`J(r)`$ on $`N`$, whose first terms are
$$J(r)=J_02\mathrm{e}^r\tau +\mathrm{e}^{2r}(2|\tau |^2J_0_T\tau )+\mathrm{},$$
and a unique *finite* jet of Kรคhler-Einstein metric $`g_{KE}`$ on $`M`$, that is locally determined up to order 2 by $`(M,\theta )`$: given some choice of coframe $`\theta ^1\mathrm{\Omega }^{1,0}H`$, the expression of its Kรคhler form $`\omega `$ is
(9)
$$\begin{array}{cc}\hfill \omega & =\mathrm{e}^r(dr\theta +d\theta )\frac{R}{2}d\theta \hfill \\ & +\frac{4}{3}\left(\frac{i}{8}R_{,\overline{1}}\vartheta ^0\theta ^{\overline{1}}\frac{i}{8}R_{,1}\vartheta ^{\overline{0}}\theta ^1\frac{1}{2}\tau _{\overline{1},1}^1\vartheta ^0\theta ^{\overline{1}}\frac{1}{2}\tau _{1,\overline{1}}^{\overline{1}}\vartheta ^{\overline{0}}\theta ^1\right)\hfill \\ & \frac{\mathrm{\Delta }_HR}{2}\mathrm{e}^rd\theta \frac{2}{3}\left(\frac{R^2}{8}|\tau |^2\frac{\mathrm{\Delta }_HR}{6}+\frac{2i}{3}(\tau _{\overline{1},11}^1\tau _{1,\overline{1}\overline{1}}^{\overline{1}})\right)\mathrm{e}^rdr\theta \hfill \\ & +\frac{2}{3}\left(\frac{R^2}{8}|\tau |^2\frac{\mathrm{\Delta }_HR}{12}\frac{i}{3}(\tau _{\overline{1},11}^1\tau _{1,\overline{1}\overline{1}}^{\overline{1}})\right)\mathrm{e}^rd\theta +o(\mathrm{e}^{2r}),\hfill \end{array}$$
where $`(\vartheta ^0=\mathrm{e}^rdr+i\theta ,\vartheta ^1)`$ is a coframe of $`\mathrm{\Omega }^{1,0}N`$ associated to $`J(r)`$.
It is explained in why higher order terms in $`\omega `$ are irrelevant in all what concerns the $`\nu `$-invariant to be defined below. We will denote by $`g_{KE}`$ the metric on $`N`$ given by this second order jet of Kรคhler metric $`g_{KE}=\omega (,J(r))`$. We observe that $`g_{KE}`$ has an *universal* polynomial expression in the powers of $`\mathrm{e}^r`$, with coefficients that are tensorial in the covariant derivatives of $`R`$ and $`\tau `$. By construction the leading term of $`g_{KE}`$ is $`g_0`$ as given in (8), and the family of metrics $`h(r)`$ induced on
$$M_r=\{r\}\times MM$$
is asymptotic to $`h_0(r)`$ in (8).
Finally, an important point here is that, although we have chosen a contact form to write down the formulas for $`g_{KE}`$, actually it does depend only on the CR structure, not on the pseudohermitian structure. This is because the filling complex structure on $`N`$ depends only on $`J`$, as does the zeroth order term of $`g_0`$, and the finite part of the Kรคhler-Einstein metric that we need is uniquely determined.
We can now define the $`\nu `$-invariant of $`M`$. According to , it is obtained by taking the limit as $`r`$ goes to infinity (i.e. by taking the diabatic limit) of the boundary contribution on $`M_r`$ of the Atiyah-Patodi-Singer formula for the characteristic number $`\chi 3\tau `$ of $`[r_0,r]\times MN`$, *with respect to the metric* $`g_{KE}`$.
###### 2.1 Definition.
The $`\nu `$-invariant of $`M`$ is
$$\nu (M)=\underset{r+\mathrm{}}{lim}B(g_{KE},M_r)3\eta (h(r)),$$
where $`\eta (h(r))`$ is the $`\eta `$-invariant of the boundary operator for the signature $`S=(1)^p(dd)`$ on $`\mathrm{\Omega }^{2p}M_r`$ (see ) with the metric $`h(r)`$, and $`B(g_{KE},M_r)`$ is an integral over $`M_r`$ of the relevant secondary characteristic class, tensorially constructed from the curvature of $`g_{KE}`$ and the second fundamental form of $`M_rN`$.
It is shown in that this limit exists and actually gives rise to a CR invariant of $`M`$ (independent on the choice of the contact form $`\theta `$). The interested reader is referred to \[11, (7.7)\] for the general formula. We will not need the precise form of the correction term $`B(g_{KE},M_r)`$ in this paper.
## 3. The renormalized eta invariant
From its very definition, the invariant $`\nu `$ is a renormalisation of $`\eta `$-invariants of a jet $`h(r)`$ of the very natural Kรคhler metric $`g_{KE}`$ restricted to slides of large radii. However, these metrics are quite intricate (as formula (9) obviously shows), and $`\nu `$ itself is given by a limit of some complicated expressions built from these metrics. For these reasons we would like to describe how $`\nu `$ is related to the $`\eta `$-invariants of the much simpler contact-rescaling family of metrics of formula (1):
$$h_\epsilon =\epsilon ^1\theta ^2+\gamma .$$
This can be done as follows: although $`\eta `$ is a priori not locally computable from the metric, its variation is. Indeed from the Atiyah-Patodi-Singer formula and Chern-Simonsโ theory one has
(10)
$$\eta (h_{\epsilon _1})\eta (h_{\epsilon _0})=\frac{1}{3}_MTp_1(_{\epsilon _1},_{\epsilon _0}),$$
where $`Tp_1(_{\epsilon _1},_{\epsilon _0})`$ is Chern-Simonsโ transgression form of the first Pontrjagin class relative to the Levi-Civita connections of the *product* metrics
$$\stackrel{~}{g}_\epsilon =dr^2+h_\epsilon \text{on}N=\times M.$$
If $`_{\epsilon _1}=_{\epsilon _0}+\alpha `$ and $`\mathrm{\Omega }_t`$ is the curvature $`2`$-form of $`_{\epsilon _0}+t\alpha `$, then
(11)
$$Tp_1(_{\epsilon _1},_{\epsilon _0})=2_0^1P_1(\alpha ,\mathrm{\Omega }_t)๐t=\frac{1}{4\pi ^2}_0^1\mathrm{Tr}(\alpha \mathrm{\Omega }_t)๐t.$$
This leads quickly to the following lemma.
###### 3.1 Lemma.
Let $`(M^3,J,\theta )`$ be a strictly pseudoconvex pseudohermitian manifold, with metric $`\gamma =d\theta (,J)`$ on the contact distribution. Then the $`\eta `$-invariants of the family of metrics $`h_\epsilon =\epsilon ^1\theta ^2+\gamma `$ have a decomposition in homogeneous terms:
(12)
$$\eta (h_\epsilon )=\underset{i=2}{\overset{2}{}}\eta _i(M,\theta )\epsilon ^i.$$
The terms $`\eta _i`$ for $`i0`$ are integral of local pseudohermitian invariants of $`(M,\theta )`$, and the $`\eta _i`$ for $`i>0`$ vanish when the torsion vanishes.
###### Proof.
Denote by $`^W`$ the Tanaka-Webster connection, with $`\tau `$ being the torsion seen as a trace-free symmetric endomorphism of $`H=\mathrm{ker}\theta `$, $`\tau ^1`$ (resp. $`\tau ^{\overline{1}}`$) being its expression as a $`(0,1)`$-form (resp $`(1,0)`$-form) relative to a choice of complex coframe $`\theta ^1`$. One computes easily the difference $`a=_\epsilon ^W`$ (see the formulas in \[44, page 316\]), and the result is a decomposition into homogeneous terms of degrees $`1`$, $`0`$ and $`1`$:
(13)
$$_\epsilon ^W=a=\underset{1}{\overset{1}{}}a^{(i)}\epsilon ^i,$$
where each $`a^{(i)}`$ is locally defined by the pseudohermitian structure: $`a^{(0)}`$ and $`a^{(1)}`$ are horizontal, but $`a^{(1)}`$ is vertical, more precisely, for horizontal $`X,YH`$ one has
$`a_X^{(1)}Y`$ $`=\gamma (\tau (X),Y)T,`$
$`a_X^{(0)}T`$ $`=\tau (X),`$
$`a_T^{(1)}Y`$ $`={\displaystyle \frac{1}{2}}JY.`$
The output is the following decomposition for the curvature
(14) $`\mathrm{\Omega }(_\epsilon )`$ $`=\mathrm{\Omega }(^W)+d^Wa+aa`$
(15) $`={\displaystyle \underset{1}{\overset{1}{}}}\mathrm{\Omega }^{(i)}\epsilon ^i.`$
Indeed, the terms $`\mathrm{\Omega }^{(\pm 2)}=a^{(\pm 1)}a^{(\pm 1)}`$ clearly vanish. Moreover,
$$\mathrm{\Omega }^{(1)}=da^{(1)}+a^{(1)}a^{(0)}+a^{(0)}a^{(1)}$$
vanishes when the torsion vanishes. From equation (11) one has
$$\epsilon \frac{d}{d\epsilon }\eta (h_\epsilon )=\frac{1}{12\pi ^2}_M\mathrm{Tr}(\mathrm{\Omega }\epsilon \frac{da}{d\epsilon })=\underset{\begin{array}{c}2i2\\ i0\end{array}}{}i\eta _i\epsilon ^i$$
where the $`\eta _i`$ ($`i0`$) are local pseudo-hermitian invariants. When the torsion vanishes, $`a^{(1)}`$ and $`\mathrm{\Omega }^{(1)}`$ vanish, so that $`\eta _i`$ vanishes for each $`i>0`$. โ
From the conformal invariance of the $`\eta `$-invariant, one deduces moreover immediately that, for a real number $`\lambda >0`$,
(16)
$$\eta _i(M,\lambda \theta )=\lambda ^i\eta _i(M,\theta ).$$
so that $`\eta _0(M,\theta )`$ is scale (but not conformally) invariant.
###### 3.2 Definition.
Let $`(M^3,\theta )`$ be a compact strictly pseudoconvex pseudohermitian manifold. The *renormalized $`\eta `$-invariant of $`(M,\theta )`$* is the constant term $`\eta _0(M,\theta )`$ in the expansion (12) for the $`\eta `$-invariants of the family of metrics $`h_\epsilon =\epsilon ^1\theta ^2+d\theta (,J)`$.
In the case where the torsion of $`(M,\theta )`$ vanishes, the terms $`\eta _i(M,\theta )`$ in (12) for $`i>0`$ vanish, so that, when $`\epsilon `$ goes to infinity instead of $`0`$, one has
(17)
$$\eta _0(M,\theta )=\underset{\epsilon \mathrm{}}{lim}\eta (h_\epsilon ):=\eta _{\mathrm{ad}}.$$
This corresponds to the geometric situation when the Reeb flow preserves the metric. Then, when $`\epsilon \mathrm{}`$, the family of metrics $`h_\epsilon `$ collapses with bounded connection and curvature. This is the well-known *adiabatic limit*, and $`\eta _0(M,\theta )`$ is then the adiabatic limit $`\eta _{\mathrm{ad}}`$ of the $`\eta `$-invariant. It has been much studied, in particular in the geometrically meaningful situation when the Riemannian flow comes from some fibration in circles over a surface .
However, we are more interested in this paper in the opposite direction: the diabatic limit, or equivalently the case where $`\epsilon `$ goes to $`0`$. Although we will not need its precise expression, making the calculations in the proof of lemma 3.1 explicit shows the term $`\eta _2(M,\theta )`$ never vanishes on contact manifolds, and has to be of type $`C_M\theta d\theta `$ for some universal non-zero constant $`C`$. Therefore $`\eta (h_\epsilon )`$ always diverges at speed $`\epsilon ^2`$ in the diabatic limit, but the constant term $`\eta _0(M,\theta )`$ is still well-defined. We called it the *renormalized* $`\eta `$-invariant, as it is reminiscent of other similar contexts where renormalized invariants have been defined .
## 4. The relation between $`\nu `$ and $`\eta _0`$
Our goal now is to prove Theorem 1.1, *i.e.* to show that on any CR manifold the $`\nu `$-invariant is related to $`\eta _0`$ in a simple way.
###### 4.1 Lemma.
There exist two constants $`C_1`$ and $`C_2`$ such that for any CR strictly pseudoconvex pseudohermitian manifold $`(M^3,J,\theta )`$, one has
(18)
$$\nu (M)+3\eta _0(M,\theta )=C_1_MR^2\theta d\theta +C_2_M|\tau |^2\theta d\theta ,$$
where $`\eta _0(M,\theta )`$ is the renormalized $`\eta `$-invariant of $`(M,\theta )`$, and $`R`$, $`\tau `$ are the Tanaka-Webster curvature and torsion of $`M`$.
One can therefore look at $`\nu (M)/3`$ as a local CR-conformal correction of $`\eta _0(M,\theta )`$ (recall that $`\eta _0(M,\theta )`$ is a priori only invariant under the rescaling $`\theta \lambda \theta `$ for $`\lambda `$ *constant*).
###### Proof.
The metrics $`g_{KE}`$ and $`h(r)=g_{KE}|_{\{r\}\times M}`$ issued from (9) are quite complicated, but are corrections of the model metrics $`g_0`$ and $`h_0(r)`$ defined in (8). More precisely, their expressions are universal polynomials in $`\mathrm{e}^r`$ and pseudohermitian invariant of $`(M,\theta )`$, and they do not actually depend on the choice of framing (except $`\theta `$) and the constants in front of each such term are universal, i.e. independent of the manifold. Therefore, using a transgression formula as in (10) and (11), but between $`h(r)`$ and $`h_0(r)`$, we see that $`\eta (h(r))\eta (h_0(r))`$ has to be an invariant universal expression of type
(19)
$$\underset{k=n}{\overset{n}{}}\mathrm{e}^{kr}_MP_k(R,\tau ,R,\tau ,\mathrm{}).$$
From lemma 3.1, and the fact that the metric $`h_0(r)`$ is $`\epsilon ^1h_\epsilon `$ with $`\epsilon =\mathrm{e}^r`$, the same holds true for $`\eta (h(r))\eta _0(M,\theta )`$.
Moreover, the boundary contribution $`B(g_{KE},M_r)`$ arising in definition 2.1 of $`\nu `$ is the integral of a secondary class built from the curvature of $`g_{KE}`$ and has therefore a development of the same type as (19). The expression
$$\nu (r)+3\eta _0(M,\theta )=B(g_{KE},M_r)3(\eta (h(r))+\eta _0(M,\theta ))$$
has then a development of the same kind. Note that this expression is void of terms in $`\mathrm{e}^{kr}`$ for $`k>0`$ since we already know from definition 2.1 and that it converges when $`r`$ goes to infinity. As a result, the local boundary contribution necessarily cancels all divergent terms, and adds (still local) convergent terms. Identifying the constant terms we get eventually:
$$\nu (M)+3\eta _0(M,\theta )=_MP_\theta (R,\tau ,R,\tau ,\mathrm{})\theta d\theta $$
where $`P_\theta `$ is some pseudohermitian local tensorial invariant. The invariance under the rescaling $`\theta \lambda ^2\theta `$ shows that the polynomial $`P_\theta `$ must satisfy
$$P_{\lambda ^2\theta }=\lambda ^4P_\theta .$$
The list of all possible expressions is easily established. Indeed, elementary invariant theory yields that such $`U(1)`$-invariant polynomials have to be sums of full contractions. Curvature $`R`$ and torsion $`\tau `$ (here we see the torsion $`\tau `$ as a tensor of type $`\tau =A_{11}\theta ^1\theta ^1`$ using some coframe $`\theta ^1`$ of $`T^{1,0}H`$) are homogeneous of weight $`2`$ with respect to the previous rescaling, while a covariant differentiation along $`T`$ decreases the weight by $`2`$, and an horizontal one by $`1`$ . Following proposition 5.13 in , we find that $`P_\theta `$ is a combination of
(20)
$$\begin{array}{c}R^2,|\tau |^2=|A_{11}|^2,R_{,0}=dR(T),\mathrm{\Delta }_HR,\\ _{0,1}^2\tau =A_{11,\overline{1}\overline{1}},_{1,0}^2\overline{\tau }=A_{\overline{1}\overline{1},11}.\end{array}$$
Full divergences do not contribute after integration over $`M`$, so that one may forget the last four expressions, and the proof of lemma 4.1 is over. โ
### Computation of the constants
We are left with the determination of $`C_1`$ and $`C_2`$ in lemma 4.1. This shall come from an explicit study of left-invariant CR structures on the three sphere.
Choose a basis $`(\alpha _1,\alpha _2,\alpha _3)`$ of left-invariant 1-forms on $`๐^3`$, such that $`d\alpha _1=\alpha _2\alpha _3`$, etc. The $`\eta `$-invariant of the left-invariant metric $`\lambda _1^2\alpha _1^2+\lambda _2^2\alpha _2^2+\lambda _3^2\alpha _3^2`$ has been computed by Hitchin \[28, formula (10)\]<sup>1</sup><sup>1</sup>1There is a slight mistake in by a factor $`2`$, as can be seen by comparing the results in for the standard sphere to those of theorem 5.2 below: one must find $`\eta _0(๐^3,std)=\frac{2}{3}`$ in the equation (22) below, rather than $`\frac{4}{3}`$ computed by .:
(21)
$$\eta (\lambda _1^2\alpha _1^2+\lambda _2^2\alpha _2^2+\lambda _3^2\alpha _3^2)=\frac{2}{3}\left(\frac{s_1^34s_1s_2}{s_3}+9\right)$$
where the $`s_i`$ are the symmetric polynomials in the $`\lambda _i^2`$. As a result, we get
$$\begin{array}{c}\eta (\alpha _1^2+\lambda _2^2\alpha _2^2+\lambda _3^2\alpha _3^2)\hfill \\ \hfill =\frac{2}{3\lambda _2^2\lambda _3^2}\left(\lambda _3^6(1+\lambda _2^2)\lambda _3^4(\lambda _2^43\lambda _2^2+1)\lambda _3^2+(\lambda _2^6\lambda _2^4\lambda _2^2+1)\right)\end{array}$$
and taking the constant term in the diabatic limit $`\lambda _3\mathrm{}`$ (*i.e.* taking $`\theta =\alpha _3`$) leads to
(22)
$$\eta _0(\alpha _1^2+\lambda ^2\alpha _2^2)=\frac{2}{3\lambda ^2}(\lambda ^4+3\lambda ^21).$$
On the other hand, the $`\nu `$-invariant can be estimated from the $`\mu `$-invariant introduced by Burns and Epstein for embeddable CR structures, or more generally CR manifolds with trivial holomorphic part of the contact bundle : for the contact form $`\theta =\alpha _3`$ and a metric $`\gamma =\lambda ^1(\alpha _1)^2+\lambda (\alpha _2)^2`$, $`\mu `$ is calculated in \[14, 4.1.A\]. Since
(23)
$$R=\frac{1+\lambda ^2}{2\lambda },|\tau |=\frac{1\lambda ^2}{2\lambda },$$
one has
$$\mu (\lambda ^1\alpha _1^2+\lambda \alpha _2^2)=\frac{1}{16\pi ^2}_{S^3}(4|\tau |^2R^2)\theta d\theta =1+\frac{3(1\lambda ^2)^2}{4\lambda ^2}.$$
It is proved in that, for a deformation of the standard CR $`3`$-sphere, one has $`\nu =3\mu +2`$, and therefore
(24)
$$\nu (\lambda ^1\alpha _1^2+\lambda \alpha _2^2)=1+\frac{9(1\lambda ^2)^2}{4\lambda ^2}.$$
From equations (22), (24) and (23) we deduce
$$(\nu +3\eta _0)(\lambda ^1\alpha _1^2+\lambda \alpha _2^2)=\frac{(1+\lambda ^2)^2}{4\lambda ^2}=\frac{1}{16\pi ^2}_{S^3}R^2\theta d\theta .$$
This yields $`16\pi ^2C_1=1`$ and $`C_2=0`$ and the proof of theorem 1.1 is done. โ
###### 4.2 Remark.
From Theorem 1.1, we see that $`3\eta _0+{\displaystyle \frac{1}{16\pi ^2}}{\displaystyle _M}R^2\theta d\theta `$ is a CR invariant. This fact can be proved directly: standard calculations in pseudohermitian geometry lead easily to the conclusion that it is invariant under conformal transformations $`\theta f\theta `$.
This provides an alternative (and independent) definition of the $`\nu `$-invariant. The latest is clearly simpler than the one explained in section 2: this is useful for computations and theoretical aspects, in particular the relation with the $`\eta `$-invariant of the contact operator $`D`$ on vertical 2-forms, as we shall see in the following sections. On the other hand, very important for the applications is the fact that $`\nu `$ arises as a boundary term in the integral of characteristic classes (see for example corollary 1.9), and this can be obtained only through the first definition and the work done in .
One may also think that this remark could serve as a basis for defining a version of $`\nu `$ in higher dimensions, by looking for local corrections of $`\eta _0`$ that would lead to a CR invariant. However, this seems a very difficult task, as the range of possible terms of the right weight is in general much larger than in (20), even in the next relevant dimension $`7`$.
## 5. Computation of the invariant on Seifert manifolds
This section is devoted to explicit computations of the $`\nu `$-invariant on $`๐^1`$-invariant CR manifolds of dimension $`3`$. Although certainly a digression from our main route towards Theorems 1.4 and 1.5, this appears as a nice direct application of the results obtained in the previous section. We have thus chosen to interrupt the pace of our proofs, and to offer this section as a refreshing intermezzo before the analytical technicalities that will follow.
We first describe our family of spherical $`3`$-dimensional compact strictly pseudoconvex CR manifolds in greater detail.
###### 5.1 Definition.
A CR-Seifert manifold is a $`3`$-dimensional compact manifold endowed with both a pseudoconvex CR structure $`(H,J)`$ and a Seifert structure, that are compatible in the following sense: the circle action $`\phi :๐^1\mathrm{Diff}(M)`$ preserves the CR structure and is generated by a Reeb field $`T`$.
Any $`๐^1`$-invariant CR structure admits a $`๐^1`$-invariant contact form $`\theta `$ if the manifold is orientable (this is proved in ). Moreover it is easily proved that that existence of a Reeb field $`T`$ (defined by $`\theta (T)=1`$ and $`\iota _Td\theta =0`$) satisfying $`\phi _{}(\frac{d}{dt})=T`$ and $`_T\theta =0`$, $`_TJ=0`$, is equivalent to the existence of a locally free action of $`๐^1`$ whose (never vanishing) infinitesimal generator preserves $`H`$ and $`J`$ and is transverse everywhere to $`H`$. Hence, our CR-Seifert manifolds could also be called *transverse $`๐^1`$-invariant CR manifolds*; note moreover that there exists a much larger class of $`๐^1`$-invariant CR manifolds, with the infinitesimal generator being sometimes tangent to the contact distribution .
As we do not assume the action to be free but only locally free, the quotient space $`\mathrm{\Sigma }=M/๐^1`$ is a surface with possibly conical singularities. Each CR-Seifert manifold is then an orbifold bundle over the compact Riemannian orbifold surface $`\mathrm{\Sigma }`$. If $`\mathrm{\Sigma }`$ is such a surface, endowed with a complex structure, orbifold $`๐^1`$-bundles are classified by their (rational) degrees $`d`$. Singularities of the bundle are located above the singularities of $`\mathrm{\Sigma }`$ in such a way that the resulting $`3`$-manifold is smooth: if the local fundamental group is $`/\alpha `$ ($`\alpha ^{}`$), a generator acts on a local chart around $`p`$ of the basis manifold as $`\mathrm{e}^{i\frac{2\rho \pi }{\alpha }}`$ and on the fiber as $`\mathrm{e}^{i\frac{2\pi \beta }{\alpha }}`$ with $`\rho `$ and $`\beta `$ prime to $`\alpha `$ (the extra parameter $`\rho `$ may seem pointless as it is always possible to reduce oneself to two parameters by taking $`\rho ^{}=1`$ and $`\beta ^{}=\beta \rho ^1`$ mod. $`\alpha `$, but this extended description will prove useful when specializing our computations to the case of lens spaces in section 10) . Any choice of equivariant connection $`1`$-form $`\theta `$ on $`M`$ endows it with an invariant CR structure, $`H`$ being chosen as the horizontal space for the connection and $`J`$ being pulled back from the base. It is strictly pseudoconvex if $`d<0`$. The interested reader is referred to for a very readable account on orbifold bundles over orbifold surfaces. Note moreover that one has
$$_M\theta d\theta =4\pi ^2d,$$
and that the metric $`\gamma =d\theta (,J)`$ projects downwards to a metric on $`\mathrm{\Sigma }`$ of volume
$$_\mathrm{\Sigma }๐\theta =2\pi d,$$
(see again for integration of forms over orbifolds). Its curvature $`R`$ equals the Tanaka-Webster curvature of $`(M,\theta )`$ and Gauss-Bonnet reads
$$_\mathrm{\Sigma }R๐\theta =2\pi \chi ,$$
where $`\chi `$ is the (rational) Euler characteristic of $`\mathrm{\Sigma }`$.
### Computations in constant curvature
In the first half of this section, we moreover assume that $`\gamma `$ has constant curvature $`R`$. In this case, the CR structure is spherical, that is $`M`$ is locally isomorphic to the standard $`3`$-sphere. Conversely, it is known that spherical CR-Seifert manifolds are exactly those of constant Tanaka-Webster curvature $`R`$, except if the base is a sphere, see for instance .
The computations now rely on the explicit derivation of the $`\eta `$-invariant of (orbifold) circle bundles over (orbifold) Riemannian surfaces with constant curvature that have been done by Komuro and more generally by Ouyang . In our conventions and notations, their results read:
###### 5.2 Theorem (Ouyang).
The $`\eta `$-invariant of the metric $`t^2\theta ^2+\gamma `$ on $`M`$ is equal to
$$\frac{1}{3}\left(d+3+2d\left(\frac{\pi t^2}{V}\chi \frac{\pi ^2t^4}{V^2}d^2\right)\right)+4\underset{j=1}{\overset{p}{}}s(\alpha _j,\rho _j,\gamma _j),$$
where $`s(\alpha ,\rho ,\gamma )=\frac{1}{4\alpha }_{k=1}^{\alpha 1}\mathrm{cot}(\frac{k\rho \pi }{\alpha })\mathrm{cot}(\frac{k\beta \pi }{\alpha })`$ is the classical Rademacher-Dedekind sum.
We can now proceed to the computation of $`\nu `$ in the constant curvature case. We have to show Corollary 1.3, which we restate here:
###### 5.3 Corollary.
Let $`M`$ be a compact $`๐^1`$-orbifold bundle of rational degree $`d<0`$ over a compact orbifold surface $`\mathrm{\Sigma }`$ of constant curvature and rational Euler characteristic $`\chi `$. Then,
(25)
$$\nu (M)=d3\frac{\chi ^2}{4d}12\underset{j=1}{\overset{p}{}}s(\alpha _j,\rho _j,\beta _j).$$
Let us remark that the $`\nu `$-invariant depends only on the topology for this class of CR manifolds, and not, for instance, on the complex structure of $`\mathrm{\Sigma }`$. This is a priori known, since the gradient of $`\nu `$ is the Cartan curvature \[11, Theorem 8.1\], which vanishes for spherical CR manifolds.
###### Proof.
According to Theorem 1.1, the $`\nu `$-invariant is given by adding a local term to the renormalized $`\eta `$-invariant. On $`๐^1`$-invariant CR manifolds with constant curvature, the renormalized invariant is easily read from Ouyangโs Theorem 5.2 above:
(26)
$$\eta _0=1+\frac{d}{3}+4\underset{j=1}{\overset{p}{}}s(\alpha _j,\rho _j,\beta _j).$$
Moreover, the integral term is just
$$\frac{1}{16\pi ^2}_MR^2\theta d\theta =\frac{4\pi ^2d\left(\frac{\chi }{d}\right)^2}{16\pi ^2}=\frac{\chi ^2}{4d},$$
which shows also Theorem 1.3 in the constant curvature case. โ
###### 5.4 Remark.
Corollary 1.3 can also be obtained by direct calculation from the original definition of $`\nu `$ and Ouyangโs formula. Indeed the asymptotically Kรคhler-Einstein metric $`g_{KE}`$ on $`[r_0,+\mathrm{}[\times M`$ can be handled with bare hands in this simple situation, and the boundary contribution counterbalancing the divergence of the sequence of $`\eta `$-invariants can be explicitly derived. Putting together Ouyangโs theorem 5.2 and these local computations yield the value of $`\nu `$, see for similar computations. This is of course a painful method, but it is still a reasonably simple case where the cancellation of divergences by local terms can be observed in detail.
### Extension to cases of non-constant curvature
We now extend the computations of $`\nu `$ to an (almost) complete proof of theorem 1.2. It is shown in that there always exist a unique (up to equivalence) transverse $`๐^1`$-contact form on an orientable Seifert manifold (careful: this might be wrong for a non-transverse action). Given the natural contact form that fixes the length of the regular fibers to $`2\pi `$, the choice of a CR structure is then equivalent to the choice of a downwards orbifold Riemannian metric $`\gamma `$ of fixed volume $`d\theta `$, and this metric might or might not be of constant curvature.
In case the base is smooth (no orbifold singularities), it is known that the adiabatic limit $`\eta _{\mathrm{ad}}`$ does not depend on the underlying metric on $`\mathrm{\Sigma }`$, see *e.g.* . As one can always find a constant curvature metric of volume $`d\theta `$ (easy consequence of Moserโs lemma on volume forms), the previous formula (26) for $`\eta _0=\eta _{ad}`$ applies. Then Theorem 1.1 enables to conclude that
(27)
$$\nu (M)=d312\underset{j=1}{\overset{p}{}}s(\alpha _j,\rho _j,\beta _j)+\frac{1}{8\pi }_\mathrm{\Sigma }R^2๐\theta .$$
If orbifolds singularities are present, it is known that every orbifold surface has a constant curvature metric, except some exceptional cases on the sphere described in . As the set of compatible complex structures with a given contact structure is contractible, this means that, except on the exceptional cases we have just alluded to, it suffices to check the following:
###### 5.5 Lemma.
The variations of $`\eta _0`$ with respect to the complex structure vanish when the torsion is zero.
###### Proof.
From Theorem 1.1, $`\eta _0`$ has the same variation as
$$\frac{\nu }{3}+\frac{1}{48\pi ^2}_MR^2\theta d\theta .$$
The variation of $`\nu `$ with respect to $`J`$ has been computed in \[11, Theorem 8.1\], namely
(28)
$$\frac{d\nu }{dJ}=\frac{3}{8\pi ^2}_MQ_J,\dot{J}\theta d\theta ,$$
where $`Q_J=iQ_{1}^{}{}_{}{}^{\overline{1}}\theta ^1Z_{\overline{1}}iQ_{\overline{1}}^{}{}_{}{}^{1}\theta ^{\overline{1}}Z_1\mathrm{End}(H)`$ is Cartanโs tensor. Its expression in term of derivatives of Tanaka-Webster curvature and torsion is given by
(29)
$$Q_{1}^{}{}_{}{}^{\overline{1}}=\frac{1}{6}R_{,1}^{}{}_{}{}^{\overline{1}}+\frac{i}{2}RA_{1}^{}{}_{}{}^{\overline{1}}A_{1}^{}{}_{}{}^{\overline{1}}{}_{,\mathrm{\hspace{0.17em}0}}{}^{}\frac{2i}{3}A_{1}^{}{}_{}{}^{\overline{1}}{}_{,\overline{1}}{}^{}{}_{}{}^{\overline{1}}.$$
On the other hand the variation of the Tanaka-Webster curvature is computed e.g in \[18, (2.20)\], and is given by
(30)
$$\dot{R}=i(E_1^{\overline{1}}{}_{,\overline{1}}{}^{}{}_{}{}^{1}E_{\overline{1}}^{1}{}_{,1}{}^{}{}_{}{}^{\overline{1}})(A_{1}^{}{}_{}{}^{\overline{1}}E_{\overline{1}}^{1}+A_{\overline{1}}^{1}E_1^{\overline{1}}),$$
where
(31)
$$\dot{J}=2E_1^{\overline{1}}\theta ^1Z_{\overline{1}}+2E_{\overline{1}}^{1}\theta ^{\overline{1}}Z_{\overline{1}}.$$
Putting everything together and integrating by parts shows that, in vanishing torsion, $`\eta _0`$ does not depend on the complex structure as needed. โ
###### 5.6 Remark.
This computations of variations may be seen as an alternative mean to determine the constant $`C_1=\frac{1}{16\pi ^2}`$ in Lemma 4.1, independently of the computations of examples done in section 4. Moreover it shows that $`\eta _0`$ is independent of $`J`$ whenever the torsion vanishes, *without any assumption on the quotient structure of $`M`$ by the Reeb flow*. This last fact will be used in section 9.
In the remaining exceptional cases over $`๐^2`$ described in , the results stay the same but the proof above does not apply anymore and one has to rely on a different technique: this will be done below in section 8.
## 6. The contact complex and the diabatic limit.
Theorem 1.1 gives a simple formula relating the $`\nu `$-invariant and the renormalized $`\eta `$-invariant $`\eta _0`$ of the contact-rescaling. According to (17), $`\eta _0`$ coincides with the adiabatic limit of $`\eta `$ in the case the CR manifold has vanishing torsion, and this enables computations, for explicit expressions of the adiabatic limit are known in a number of cases. But a deeper question is to relate directly the $`\nu `$-invariant to the geometry and spectral theory of the CR or pseudohermitian manifold.
In the sequel we shall consider a natural $`\eta `$-invariant arising in pseudohermitian geometry. One actually knows by a candidate for this, coming from the contact-de Rham complex. We shall briefly recall its construction in dimension $`3`$ and its relation with the diabatic limit.
Let $`M`$ be a $`3`$-dimensional contact manifold and $`H`$ its contact distribution. We denote by $`\mathrm{\Omega }^{}H`$ the space of horizontal forms, *i.e.* the space of sections of the alternating algebra over the dual of the bundle $`H`$. Let also $`\mathrm{\Omega }^{}V`$ be the subspace of vertical forms on $`M`$, by which we mean โtrueโ forms in $`\mathrm{\Omega }^{}M`$ vanishing on $`H`$. Equivalently, one has $`\mathrm{\Omega }^{}V=\{\theta \alpha \}=\theta \mathrm{\Omega }^{}H`$ for any local choice of contact form $`\theta `$. The contact-de Rham complex is then the following:
(32)
$$C^{\mathrm{}}(M)\stackrel{d_H}{}\mathrm{\Omega }^1H\stackrel{D}{}\mathrm{\Omega }^2V\stackrel{d_H}{}\mathrm{\Omega }^3M,$$
where for $`fC^{\mathrm{}}(M)`$, $`d_Hf\mathrm{\Omega }^1H`$ stands for the restriction of $`df`$ to $`H`$, while
$$d_H:\mathrm{\Omega }^2V\mathrm{\Omega }^3M$$
is just de Rhamโs differential restricted to $`\mathrm{\Omega }^2V`$ in $`\mathrm{\Omega }^2M`$, and $`D`$ is defined as follows: since $`d`$ induces an isomorphism
$$d_0:\mathrm{\Omega }^1V\mathrm{\Omega }^2H\text{ with }d_0(f\theta )=fd\theta _{\mathrm{\Lambda }^2H},$$
then any $`\alpha `$ in $`\mathrm{\Omega }^1H`$ admits a unique extension $`\mathrm{}(\alpha )`$ in $`\mathrm{\Omega }^1M`$ such that $`d\mathrm{}(\alpha )`$ belongs to $`\mathrm{\Omega }^2V`$; namely, given any initial extension $`\overline{\alpha }`$ of $`\alpha `$, one has
(33)
$$\mathrm{}(\alpha )=\overline{\alpha }d_0^1(d\overline{\alpha })_{\mathrm{\Lambda }^2H}.$$
We then define
(34)
$$D\alpha =d\mathrm{}(\alpha ).$$
This differential $`D`$ is a second order operator, since the lifting $`\mathrm{}:\mathrm{\Omega }^1H\mathrm{\Omega }^1M`$ is a first order one. Moreover one sees easily that $`\mathrm{}`$ induces an homotopy equivalence between the contact and de Rham complexes, together with the natural restrictions, and the retraction $`\mathrm{}^{}:\mathrm{\Omega }^2M\mathrm{\Omega }^2V`$ defined by
$$\mathrm{}^{}(\alpha )=\alpha dd_0^1\alpha _{\mathrm{\Lambda }^2H}.$$
From now on we will suppose moreover that the contact manifold $`M`$ is endowed with a strictly pseudoconvex CR structure $`J`$, together with some choice of contact form $`\theta `$. We consider the contact-rescaling sequence of metrics of (8)
$$h_0(r)=\mathrm{e}^{2r}\theta ^2+\mathrm{e}^rd\theta (,J).$$
Let $`\epsilon =\mathrm{e}^r`$, as before, and define
(35)
$$g_\epsilon =\epsilon ^2\theta ^2+\epsilon ^1d\theta (,J)=h_0(r).$$
This metric induces an orthogonal splitting $`TM=HT`$ where $`T`$ is the Reeb field of $`\theta `$, and one can identifies $`\mathrm{\Omega }^1H`$ with โtrueโ $`1`$-forms on $`M`$ vanishing on $`T`$. Observing that Hodge $``$-operator exchanges $`\mathrm{\Omega }^1H`$ and $`\mathrm{\Omega }^2V`$ and one can consider $`D`$ acting on closed vertical $`2`$-forms $`\mathrm{\Omega }_D^2V=\mathrm{\Omega }^2V\mathrm{im}D`$.
Following \[2, Theorem 4.14\], we define the boundary operator for the signature attached to the Riemannian metric $`g_\epsilon `$ as
$$S_\epsilon =(1)^p(_\epsilon dd_\epsilon ),$$
acting on $`\mathrm{\Omega }^{2p}M=C^{\mathrm{}}M\mathrm{\Omega }^2M`$. As observed in \[2, Prop 4.20\], one may remove some spectral symmetry, and its $`\eta `$-function
(36)
$$\eta (S_\epsilon )(s)=\mathrm{Tr}^{}(S_\epsilon |S_\epsilon |^{(s+1)})=\underset{\lambda _i\mathrm{spec}(S_\epsilon )\{0\}}{}\frac{\lambda _i}{|\lambda _i|^{s+1}}$$
actually coincides with that of $`d_\epsilon `$ when restricted to $`\mathrm{\Omega }_d^2M=\mathrm{\Omega }^2M\mathrm{im}d`$. Note that we have used $`\mathrm{Tr}^{}`$ to denote a trace taken outside the $`0`$-eigenspace. In the same vein, the notation $`\mathrm{spec}^{}`$ used below will denote a spectrum *where the $`0`$-eigenvalue has been removed*.
From \[4, p. 74\] or \[24, Chap. 1.10\], the series (36) is absolutely convergent for $`\mathrm{Re}s>3`$ and has a meromorphic extension to $``$, with possibly simple poles at $`s=3n`$, $`n`$. By Atiyah-Patodi-Singerโs theorem , $`\eta (S_\epsilon )(s)`$ is actually regular at $`s=0`$ and its value there is called the $`\eta `$-invariant of $`(M,g_\epsilon )`$. Similarly, an $`\eta `$-function and its value at $`0`$ can be defined for the operator $`D`$ in dimension $`3`$. This mainly follows by applying the same ideas, but with the adequate symbolic calculus for hypoelliptic operators, see section 9.
In order to compare them, let us now compute $`d_\epsilon `$ and $`D_\epsilon `$ using the decomposition of $`\mathrm{\Omega }^2M`$ into vertical and horizontal $`2`$-forms:
$$\alpha =\theta \alpha _T+\alpha _H,$$
with $`\alpha _T\mathrm{\Omega }^1H`$, $`\alpha _H\mathrm{\Omega }^2H`$. From (35) one sees that
$$_\epsilon \alpha =\epsilon _H\alpha _T+\theta _H\alpha _H$$
where $`_H`$ denotes the induced Hodge duality on $`H`$. In matrix form, one gets
(37)
$$d_\epsilon =\left(\begin{array}{cc}\epsilon _T_H& d_H_H\\ \epsilon d_H_H& 1\end{array}\right),$$
where $`_T`$ is the Lie derivative along $`T`$.
Using (33) and (34) one finds that $`\mathrm{}(\beta )=\beta (_Hd_H\beta )\theta `$ on $`\mathrm{\Omega }^1H`$, so that $`D\beta =\theta (_T+d_H_Hd_H)\beta `$, and hence
(38)
$$D_\epsilon (\theta \alpha _T)=\epsilon \theta (_T+d_H_Hd_H)_H\alpha _T$$
on $`\mathrm{\Omega }^2V=\theta \mathrm{\Omega }^1H`$.
The whole spectrum of $`D_\epsilon =\epsilon D_1`$ then collapses at speed $`\epsilon `$ in the diabatic limit $`\epsilon 0`$, whereas part of the spectrum of $`d_\epsilon `$ is not collapsing: for instance $`(d_\epsilon )(d\theta )=d\theta `$. Hence the diabatic behaviour of the *whole* spectrum of $`d_\epsilon `$ cannot be related to $`D_\epsilon `$ alone, and indeed only the collapsing spectra are related. This shows up in the following formulas, which are direct consequences of (37) and (38), or even more directly from the definitions (33) and (34) of $`\mathrm{}`$ and $`D`$. If $`P_\epsilon =\epsilon ^1d_\epsilon `$,
(39)
$$\begin{array}{cc}\hfill P_\epsilon =\epsilon ^1d_\epsilon & =\left(\begin{array}{cc}D_1& 0\\ 0& 0\end{array}\right)+\left(\begin{array}{cc}(d_H_H)^2& \epsilon ^1d_H_H\\ d_H_H& \epsilon ^1\end{array}\right)\hfill \\ & =\mathrm{\Pi }_{\mathrm{\Omega }^2V}(D_1)\mathrm{\Pi }_{\mathrm{\Omega }^2V}+\epsilon P_\epsilon \mathrm{\Pi }_{\mathrm{\Omega }^2H}P_\epsilon .\hfill \end{array}$$
It follows that in the diabatic limit $`\epsilon 0`$ all the eventually bounded spectrum of $`P_\epsilon =\epsilon ^1d_\epsilon `$ converges, at least weakly, towards the spectrum of $`D_1`$. Actually its turns out that this spectral convergence is uniform over bounded intervals, as a consequence of the uniform convergence in the diabatic limit of the resolvents $`(\lambda P_\epsilon )^1`$ on $`\mathrm{ker}d`$ towards $`(\lambda D_1)^1`$, for $`\lambda `$ \[45, theorem 3.6\].
Such a spectral convergence is unfortunately only a first step in the study of a global spectral invariant like $`\eta `$. To illustrate this, recall that by an equivalent expression of the Riemannian $`\eta `$-invariant is given by
(40)
$$\eta (P_\epsilon )(0)=\pi ^{1/2}_0^{\mathrm{}}\mathrm{Tr}\left(P_\epsilon \mathrm{e}^{tP_\epsilon ^2}\right)\frac{dt}{\sqrt{t}}.$$
Now by \[45, Theorem 7.1\] the following global trace convergence holds
$$\mathrm{Tr}(P_\epsilon \mathrm{e}^{tP_\epsilon ^2})\mathrm{Tr}(D\mathrm{e}^{tDD^{}}),$$
when $`\epsilon `$ goes to $`0`$, but uniformly on $`t`$ *only for $`tt_0>0`$*. It cannot be true for small $`t`$ since the $`\eta `$-invariants and the integrals (40) diverge in the diabatic limit (although one knows by transgression formulas that these divergences of $`\eta (P_\epsilon )(0)`$ are given by local expressions). From the analytic viewpoint, these divergences are rooted in the transition from elliptic towards hypoelliptic operators, that cannot be uniform in all $`(t,\epsilon )`$ regimes. For instance, the asymptotic spectral densities (Weylโs laws), or the powers of $`t`$ occurring in the asymptotic expansions of the heat kernels for $`t0`$ are not the same for the elliptic $`P_\epsilon `$ and the hypoelliptic $`D`$. However it is possible, as is usual in such asymptotic spectral problems, that the divergences occurring in the $`(d_\epsilon ,D)`$ transition when $`\epsilon `$ and $`t`$ go to $`0`$, are ruled again by local expressions in the curvature, see also Remark 8.5. This would provide directly a relation like (7) between the finite part $`\eta _0`$ of $`\eta (P_\epsilon )`$ in the diabatic limit and the contact $`\eta `$-invariant $`\eta (D)`$. Unfortunately, the techniques used in cannot handle these problems in the general case. The analysis can however be done in the particular case of CR-Seifert manifolds, and we will now restrict ourselves to this case.
## 7. Spectral analysis on Seifert manifolds.
As explained above, we will now deal with CR-manifolds endowed with both a Seifert and a CR structure compatible in the sense that the circle action $`\phi :๐^1\mathrm{Diff}(M)`$ preserves the CR structure $`(H,J)`$ and is generated by a Reeb field $`T`$. An invariant contact form $`\theta `$ has then been chosen, and we note that in this section, opposite to section 5, we will never assume the Webster curvature to be constant.
The circle action allows to perform a Fourier decomposition of functions or forms, inside $`M`$ and without referring to the quotient structure. For instance, given $`n`$ and $`fC^0(M)`$, its $`n`$-th component is the function on $`M`$ defined by
$$\pi _nf=\frac{1}{2\pi }_0^{2\pi }\mathrm{e}^{int}(f\phi _t)๐t.$$
It satisfies $`(\pi _nf)\phi _t=\mathrm{e}^{int}(\pi _nf)`$, so that $`_T(\pi _nf)=in\pi _nf`$ on $`C^1(M)`$. The projections $`\pi _n`$ preserve and are clearly bounded on all $`C^p(M)`$, $`L^p(M)`$ or Sobolev spaces. Moreover, the Hilbert sum of all $`\pi _n`$ for $`n`$ in $``$ is the identity on $`L^2(M)`$. Last, this circle action preserves all structures and operators related to the above choice of contact form, so that we will be able to split their spectra into Fourier components.
We can now study the spectral aspects of the contact rescaling $`g_\epsilon `$ in (35) on a CR-Seifert manifold $`M`$. Of course the adiabatic limit exists in this situation, and has already been much studied, see *e.g.* , but we will need a different approach here, focusing on the diabatic behaviour of $`d_\epsilon `$ and $`\eta (d_\epsilon )`$, as related to the spectrum of $`D`$ and its $`\eta `$-invariant.
One computes easily the Laplacian on $`\mathrm{\Omega }^2M`$, relatively to the splitting
$$\mathrm{\Omega }^2M=\theta \mathrm{\Omega }^1H\mathrm{\Omega }^2H,$$
namely
(41)
$$\mathrm{\Delta }_\epsilon =\left(\begin{array}{cc}\epsilon \mathrm{\Delta }_H\epsilon ^2T^2& d_H_H\\ \epsilon d_H_H& 1+\epsilon \mathrm{\Delta }_H\epsilon ^2T^2\end{array}\right),$$
where $`\mathrm{\Delta }_H=d_H\delta _H+\delta _Hd_H`$ is the horizontal Laplacian (not to be confused with the contact Laplacian introduced in ), $`T`$ denotes here the Lie derivative along $`T`$, and we have used that $`T^{}=T`$ and $`[T,\delta _H]=0`$ since $`T`$ is a Killing Reeb field on the CR-Seifert manifold. We observe from (37) that the non diagonal part of $`\mathrm{\Delta }_\epsilon `$ is the same as that of $`d_\epsilon `$, so that
$$\mathrm{\Delta }_\epsilon =d_\epsilon +\epsilon \left(\begin{array}{cc}\mathrm{\Delta }_H+T_H& 0\\ 0& \mathrm{\Delta }_H\end{array}\right)\epsilon ^2T^2.$$
When studying spectral asymmetry, we restrict ourselves to the subspace $`\mathrm{\Omega }_d^2M=\mathrm{im}d`$ of $`\mathrm{\Omega }^2M`$, on which $`\mathrm{\Delta }_\epsilon =(d_\epsilon )^2`$. We get therefore the following expression relating pairwise commuting operators:
(42)
$$(d_\epsilon )^2=(d_\epsilon )+(\epsilon K)\epsilon ^2T^2,$$
with
$$K=\left(\begin{array}{cc}\mathrm{\Delta }_H+T_H& 0\\ 0& \mathrm{\Delta }_H\end{array}\right).$$
Therefore if $`\alpha \mathrm{\Omega }_d^2M\{0\}`$ satisfies
(43)
$$(d_\epsilon )\alpha =\lambda _\epsilon \alpha ,K\alpha =k\alpha \text{and}T^2\alpha =n^2\alpha ,$$
for $`\lambda _\epsilon `$ a non-zero eigenvalue of $`d_\epsilon `$, then
(44)
$$\lambda _\epsilon +\epsilon k+\epsilon ^2n^2=\lambda _\epsilon ^20,$$
and, necessarily,
(45)
$$\lambda _\epsilon =\lambda _\epsilon ^+\mathrm{or}\lambda _\epsilon ^{}\mathrm{with}\lambda _\epsilon ^\pm =\frac{1\pm \sqrt{1+4\epsilon (k+4\epsilon n^2)}}{2}.$$
Hence the spectrum of $`d_\epsilon `$ splits in two families which behave differently in the diabatic limit $`\epsilon 0`$. Eigenvalues of type $`\lambda _\epsilon ^{}`$ all collapse, while those of type $`\lambda _\epsilon ^+`$ all converge to $`1`$. According to the general results of discussed in section 6, only eigenvalues of type $`\lambda _\epsilon ^{}`$ are related to $`D`$, after rescaling by $`\epsilon ^1`$.
The previous eigenvalue equation (45) is only a necessary condition and we have to determine which of the possible $`\lambda _\epsilon ^\pm `$ are effectively present in $`\mathrm{spec}(d_\epsilon )`$ and to compute their multiplicities. To do this, we use the splitting the induced by the choice of the Reeb field: suppose $`\alpha =\theta \alpha _T+\alpha _H`$ is a $`2`$-form in the image of $`d`$. By (37), the system $`(d_\epsilon )\alpha =\lambda _\epsilon \alpha `$ is
(46) $`(\lambda _\epsilon \epsilon T_H)\alpha _T`$ $`=d_H_H\alpha _H`$
(47) $`(\lambda _\epsilon 1)\alpha _H`$ $`=\epsilon d_H_H\alpha _T.`$
Suppose now that
(48)
$$(d_\epsilon )\alpha =\lambda _\epsilon \alpha ,K\alpha =k\alpha \text{and}T^2\alpha =n^2\alpha .$$
Then we observe that $`_H=J`$ on $`\mathrm{\Omega }^1H`$ and $`(T_H)^2=T^2=n^2`$. Therefore (46) gives
(49)
$$(\lambda _\epsilon ^2\epsilon ^2n^2)\alpha _T=(\lambda _\epsilon +\epsilon T_H)d_H_H\alpha _H,$$
so that $`\alpha _H`$ determines uniquely $`\alpha _T`$ when $`\lambda _\epsilon ^2\epsilon ^2n^2`$. A first (quite large) part of the non-zero spectrum is then handled as follows.
###### 7.1 Proposition.
$``$ Forms $`\alpha =\theta \alpha _T+\alpha _H`$ in $`\mathrm{\Omega }_d^2M`$ satisfying
(50)
$$(d_\epsilon )\alpha =\lambda _\epsilon ^+\alpha ,K\alpha =k\alpha \text{and}T^2\alpha =n^2\alpha $$
are in one-to-one linear correspondence with forms $`\alpha _H`$ in $`\mathrm{\Omega }^2H`$ satisfying
(51)
$$\mathrm{\Delta }_H\alpha _H=k\alpha _H\text{and}T^2\alpha _H=n^2\alpha _H.$$
$``$ Forms $`\alpha =\theta \alpha _T+\alpha _H`$ in $`\mathrm{\Omega }_d^2M`$ satisfying
(52)
$$(d_\epsilon )\alpha =\lambda _\epsilon ^{}\alpha ,K\alpha =k\alpha \text{and}T^2\alpha =n^2\alpha $$
such that $`(\lambda _\epsilon ^{})^2\epsilon ^2n^2`$ are in one-to-one linear correspondence with forms $`\alpha _H`$ in $`\mathrm{\Omega }^2H`$ satisfying
(53)
$$\mathrm{\Delta }_H\alpha _H=k\alpha _H\text{and}T^2\alpha _H=n^2\alpha _H$$
with $`k|n|`$.
###### Proof.
Note first that, for any eigenvector $`\alpha `$ of $`d_\epsilon `$ satisfying either (50) or (52), one may have $`(\lambda _\epsilon )^2=\epsilon ^2n^2`$ only if (52) holds. Hence, in the positive case, one always has $`\alpha _H0`$, and, as a result, $`\mathrm{\Delta }_H\alpha _H=k\alpha _H`$, $`k`$ is necessarily non-negative and $`T^2\alpha _H=n^2\alpha _H`$. In the negative case, the same holds only if $`(\lambda _\epsilon ^{})^2\epsilon ^2n^2`$, and (44) shows that this is equivalent to $`k|n|`$.
Conversely, suppose now given $`\alpha _H`$, $`n`$, $`k`$, $`\lambda _\epsilon `$ as needed. From (49), one defines
$$\alpha _T=(\lambda _\epsilon ^2\epsilon ^2n^2)^1(\lambda _\epsilon +\epsilon T_H)d_H_H\alpha _H,$$
which satisfies (46). To check (47), recall that
$$\delta _H=_Hd_H_H\text{and}d_H^2=LT=TL,$$
(the last equation being a consequence of $`d^2=0`$ see e.g. \[45, p. 415\] with $`L(f)=fd\theta `$). One finds
$`(\lambda _\epsilon ^2\epsilon ^2n^2)d_H_H\alpha _T`$ $`=(\lambda _\epsilon d_H\delta _H\alpha _H\epsilon d_H^2T_H\alpha _H)`$
$`=(\lambda _\epsilon \mathrm{\Delta }_H\epsilon T^2)\alpha _H`$
$`=(\lambda _\epsilon k+\epsilon n^2)\alpha _H.`$
The eigenvalue equation (44) then easily leads to (47). โ
For later use, note that the choice $`(k,n)=(0,0)`$ in the positive case leads to $`\alpha _H=Cd\theta `$ and $`\lambda _\epsilon =1`$, hence $`\alpha _T=0`$ by (46), and this is the only case where this might happen by (47).
Proposition 7.1 shows a large part of $`\mathrm{spec}^{}(d_\epsilon )`$ is symmetric with respect to $`\frac{1}{2}`$ and is parametrised trough (45) by the spectrum $`\{k+\epsilon n^2\}`$ of the non-negative elliptic Laplacian $`L_{\epsilon ,H}=\mathrm{\Delta }_H\epsilon T^2`$ acting on $`\mathrm{\Omega }^2H`$, or equivalently by the spectrum of
(54)
$$\mathrm{\Delta }_\epsilon =\mathrm{\Delta }_H\epsilon T^2$$
acting on functions. However there are โholesโ in this symmetry corresponding to the eigenvalues $`\lambda _\epsilon ^{}=\epsilon k`$ when $`k=|n|`$, for in this case $`\alpha _T`$ is not uniquely determined by $`\alpha _H`$ so that we will have to treat these on a separate footing. This means that in the case $`\lambda _\epsilon =\lambda _\epsilon ^{}`$, we have to remove from the parameter space the horizontal forms $`\alpha _H`$ in
(55)
$$^0=\mathrm{ker}(\mathrm{\Delta }_H^2+T^2).$$
This space has a simple description using the complex structure $`J`$ and the associated splitting $`\mathrm{\Omega }^1H=\mathrm{\Omega }^{1,0}H\mathrm{\Omega }^{0,1}H`$. We recall that the component $`d_H^{0,1}`$ of $`d_H`$ from functions to $`\mathrm{\Omega }^{0,1}H`$ is called the $`\overline{}_b`$ operator, and its kernel is the space of CR functions.
###### 7.2 Proposition.
The space $`_H^0`$ is the space of pluri-CR functions, *i.e.* real parts of CR functions.
###### Proof.
Consider the Kohn Laplacians $`\overline{\mathrm{}}_b=\overline{}_b^{}\overline{}_b`$ and $`\mathrm{}_b=_b^{}_b`$ acting on functions. Following, say, \[34, Theorem 2.3\], one has in dimension $`3`$
(56)
$$\mathrm{\Delta }_H=\overline{\mathrm{}}_b+\mathrm{}_b\mathrm{and}iT=\overline{\mathrm{}}_b\mathrm{}_b.$$
Since $`T`$ commutes with everything here one gets
$$\mathrm{\Delta }_H^2+T^2=4\overline{\mathrm{}}_b\mathrm{}_b=4\mathrm{}_b\overline{\mathrm{}}_b.$$
If $`f`$ is a real function in $`^0`$ then $`g=\mathrm{}_bf`$ is CR since its image by $`\overline{\mathrm{}}_b`$ is zero, and is in the image of $`\mathrm{\Delta }_H`$ since its integral vanishes. Hence
$$\mathrm{\Delta }_Hf=\overline{\mathrm{}}_bf+\mathrm{}_bf=\overline{g}+g=2\mathrm{Re}g,$$
and $`f=2\mathrm{Re}h`$ with $`h=\mathrm{\Delta }_H^1g`$ is a CR function as needed. โ
We now study the missing case $`\lambda _\epsilon ^2=\epsilon ^2n^2`$. We first recall that complex vertical forms $`\mathrm{\Omega }^{}V\theta \mathrm{\Omega }^{}H`$ also have a natural bigrading inherited from $`J`$ on $`H`$, independently from $`\theta `$. Of particular interest here is the
###### 7.3 Definition.
The bundle $`K_M\theta \mathrm{\Omega }^{1,0}H`$ of $`2`$-forms vanishing on $`H^{0,1}`$ is called the canonical CR bundle. We denote by $`^{2,0}`$ its subspace of closed sections, also called holomorphic $`(2,0)`$-forms, and $`_+^2`$ the real part of $`^{2,0}`$.
When the CR manifold $`M`$ can be locally embedded in a $`4`$-dimensional complex manifold $`N`$, $`K_M`$ is the restriction to $`M`$ of the canonical bundle $`K_N=\mathrm{\Omega }^{2,0}N`$ of $`N`$, and holomorphic forms are local restrictions of holomorphic $`(2,0)`$-forms in $`N`$, see for instance. This explains the notation in the previous definition, as $`^{2,0}`$ (resp. $`_+^2`$) is related to the space of holomorphic $`(2,0)`$-forms in the usual sense on $`N`$ (resp. to the space of self-dual $`2`$-forms, orthogonal to the Kรคhler form). Note that this is indeed the case for our CR-Seifert manifolds for one can take $`N=M\times `$ with the extension of $`J`$ considered above.
We now show that the remaining spectrum of $`d_\epsilon `$ is entirely given by holomorphic forms.
###### 7.4 Proposition.
A $`2`$-form $`\alpha \mathrm{\Omega }_d^2M`$ satisfies
(57)
$$(d_\epsilon )\alpha =\lambda _\epsilon ^{}\alpha ,K\alpha =k\alpha \text{and}T^2\alpha =n^2\alpha $$
with $`(\lambda _\epsilon ^{})^2=\epsilon ^2n^2`$ (*i.e.* $`k=|n|`$) if and only if $`\alpha _H=0`$ and $`\alpha =\theta \alpha _T`$ belongs to $`_+^2`$.
###### Proof.
Let $`\alpha =\theta \alpha _T+\alpha _H`$ in $`\mathrm{\Omega }_d^2M`$ be an eigenfunction of $`d_\epsilon `$ satisfying (57) and $`\lambda _\epsilon ^2=\epsilon ^2n^2`$. By (44) one has also $`\lambda _\epsilon =\epsilon k`$. Since $`(T_H)^2=T^2=n^2=k^2`$ on $`\mathrm{\Omega }^1H`$, one can split
$$\alpha _T=\alpha _T^++\alpha _T^{}\mathrm{with}(T_H)\alpha _T^\pm =\pm k\alpha _T^\pm .$$
Then (46) is equivalent to
(58)
$$2\epsilon k\alpha _T^+=d_H_H\alpha _H.$$
Moreover $`K\alpha =k\alpha `$ gives $`(\mathrm{\Delta }_H+T_H)\alpha _T=k\alpha _T`$, which implies $`\mathrm{\Delta }_H\alpha _T^+=0`$ since $`[\mathrm{\Delta }_H,T_H]=0`$ on $`\mathrm{\Omega }^1H`$. Therefore $`\alpha _T^+`$ lives in $`\mathrm{ker}\delta _H`$ leading by (58) to $`\mathrm{\Delta }_H_H\alpha _H=0`$, hence to $`\alpha _H=Cd\theta `$ and $`k=n=0`$. If $`C0`$, this implies by the eigenvalue identity (44) that either $`\lambda _\epsilon =\lambda _\epsilon ^{}=0`$, which is impossible since we consider the non-zero spectrum, or to $`\lambda _\epsilon =\lambda _\epsilon ^+=1`$ which is impossible, too, because one would have $`(\lambda _\epsilon )^2\epsilon ^2n^2`$. We get then that in any case considered in the present proof, $`\alpha _H=0`$, so that $`\alpha `$ is a vertical form.
Now (47) reads $`\delta _H\alpha _T=0`$, or equivalently
$$d_H(\theta J\alpha _T)=0.$$
Recall now that $`\alpha `$ belongs to $`\mathrm{\Omega }_d^2M`$, hence is closed. The $`(1,0)`$-part of $`\alpha _T`$ is then closed and $`\theta \alpha _T`$ lives in $`_+^2`$ as needed.
Conversely, $`_+^2`$ is preserved by $`J`$ and $`T`$. Thus it can be split in eigenspaces of $`T_H=JT=k`$, on which $`d_\epsilon =k`$ by definition, see (37). โ
We now summarize our spectral study of $`d_\epsilon `$ in relation to the diabatic limit $`\epsilon 0`$.
###### 7.5 Corollary.
The spectrum of $`d_\epsilon `$ splits into the following families:
1. A converging part $`\mathrm{\Lambda }_\epsilon ^+`$, converging to $`1`$ and parametrised by the whole spectrum of $`\mathrm{\Delta }_\epsilon =\mathrm{\Delta }_H\epsilon T^2`$ (acting on functions) by the formula
$$\mathrm{\Lambda }_\epsilon ^+=\mathrm{spec}\left(\frac{1+\sqrt{1+4\epsilon \mathrm{\Delta }_\epsilon }}{2}\right).$$
2. A collapsing part, converging to $`0`$, itself divided into two families:
1. the first one $`\mathrm{\Lambda }_\epsilon ^{}`$, nearly symmetric to $`\mathrm{\Lambda }_\epsilon ^+`$:
$$\mathrm{\Lambda }_\epsilon ^{}=\mathrm{spec}\left(\frac{1\sqrt{1+4\epsilon \mathrm{\Delta }_\epsilon }}{2}\right),$$
but $`\mathrm{\Delta }_\epsilon `$ has here to be restricted to the orthogonal of the space of pluri-CR functions $`^0`$.
2. the spectrum $`\mathrm{\Lambda }_\epsilon ^0`$ of $`\epsilon T_H=\epsilon JT`$ acting on $`_+^2`$, the real parts of holomorphic forms in the canonical CR bundle.
The signs of the eigenvalues in the first two families are clear. About the third one, we can notice:
###### 7.6 Proposition.
Up to some finite dimensional space, $`d_\epsilon `$ is *positive* on $`_+^2`$.
###### Proof.
Recall that $`d_\epsilon =JT`$ on $`_+^2`$. Consider then the splitting of the Tanaka-Webster connection $`_H=_{1,0}+_{0,1}`$ on $`H`$. Then on $`K_M=\theta \mathrm{\Omega }^{1,0}H`$ one has in dimension $`3`$,
$$R=_{0,1}^{}_{0,1}_{1,0}^{}_{1,0}i_T.$$
On holomorphic forms $`^{2,0}`$ in $`K_M`$, the Lie derivative in $`T`$ equals $`_T`$ and the previous equation reduces to
$$iT=R+_{1,0}^{}_{1,0},$$
which implies that $`(iT+R)`$ is a non-negative operator. As the spectrum of $`d_\epsilon `$ (on closed forms) is discrete and without accumulation points, there is only a finite dimensional space of eigenvectors with nonpositive eigenvalues. โ
In order to get more symmetry in the spectral decomposition of $`d_\epsilon `$, one can fill in the holes in $`\mathrm{\Lambda }_\epsilon ^{}`$ by adding $`\mathrm{\Delta }_\epsilon `$ on $`_0`$. As already discussed, this corresponds to adding the cases $`k=|n|`$ and $`\lambda _\epsilon =\epsilon k0`$. Given $`k`$, the multiplicity of each added *virtual* eigenvalue $`\epsilon k`$ is equal to $`2h_0(k)`$ by Proposition 7.2, where we have denoted
$$h_0(k)=\mathrm{dim}_{}\left\{\text{CR functions }f\text{ such that }iTf=kf\right\}.$$
Observe that by (56), $`h_0(k)=0`$ if $`k<0`$. In the same spirit, the holomorphic part $`\mathrm{\Lambda }_\epsilon ^0`$ above consists in $`\{\epsilon kk^{}\}`$, with multiplicity $`2h_2(k)`$ given by
$$h_2(k)=\mathrm{dim}_{}\left\{\text{holomorphic }(2,0)\text{-forms }\alpha ^{2,0}\text{ such that }iT\alpha =k\alpha \right\}.$$
Considering the positive operators
$$Q_\epsilon ^\pm =\frac{\pm 1+\sqrt{1+4\epsilon \mathrm{\Delta }_\epsilon }}{2\epsilon },$$
leads to the more suggestive decomposition:
(59)
$$\mathrm{spec}^{}\left(\frac{d_\epsilon }{\epsilon }\right)=\pm \mathrm{spec}^{}\left(Q_\epsilon ^\pm \right)2\times \mathrm{spec}^{}\left(iT_{^{2,0}}\right)2\times \mathrm{spec}^{}\left(iT_{\mathrm{ker}\overline{}_b}\right).$$
This formula shows that the virtual spectrum of $`d_\epsilon `$ consists in a two completely different parts: a (nearly) symmetric part to $`1/2`$, that varies with $`\epsilon `$, and a *constant* holomorphic part. We will see in Lemma 8.4 that the symmetric part always contributes to $`1`$ in the renormalized $`\eta `$-invariant $`\eta _0`$ when torsion vanishes. Hence the computation of $`\eta _0`$ finally reduces to counting holomorphic objects, as will be done in section 8. This phenomenon has already been observed on a smooth base in and over orbifolds, in the adiabatic context and constant curvature, in .
## 8. The spectrum of $`D`$ and comparison of the $`\eta `$-invariants
Our goal is now to relate our description of the spectrum of $`P_\epsilon =\epsilon ^1d_\epsilon `$ to the spectrum of the middle operator of the contact complex $`D`$. We already know (see the discussion at the end of section 6) that the bounded spectrum of $`P_\epsilon `$ converges towards that of $`D`$ in the diabatic limit . Therefore from Corollary 7.5 the non-zero spectrum of $`D`$ has to split as follows
(60)
$$\mathrm{spec}^{}(D)=\mathrm{spec}^{}(\mathrm{\Delta }_H\text{on}(^0)^{})\mathrm{spec}^{}(JT\text{on}_+^2)$$
(note the lack of uniformity already noted in the introduction in the convergence of $`\mathrm{\Lambda }_\epsilon ^{}`$ when $`\epsilon 0`$, as each eigenvalue $`\mu `$ in the spectrum of $`\mathrm{\Delta }_H`$ is approached at a speed approximately $`\epsilon \mu `$). This is enough to compare the needed $`\eta `$-invariant to $`\eta _0`$ and conclude (see (64) below and the discussion following it), but we would like first to spend a few lines to reinterpret this more precisely in the CR-Seifert context.
### The spectrum of $`D`$ from the CR viewpoint
First of all, the second spectral family of eigenvalues in (60) is clearly embedded in $`\mathrm{spec}^{}(D)`$, as (38) shows that $`D=TJ`$ on $`_+^2`$. To understand where the first one comes from, we consider the following operator
$$Q=d_HJ:\mathrm{ker}d_H\mathrm{\Omega }^2V\mathrm{\Omega }^3M.$$
By definition $`_+^2=\mathrm{ker}Q`$. We also remark that
$$(Q^{})_M=(\mathrm{\Pi }_{\mathrm{ker}d_H}J\delta _H)_M=_M(\mathrm{\Pi }_{\mathrm{ker}d_H}Jd_H)$$
so that $`\mathrm{ker}Q^{}=_M^0`$ and $`\overline{\mathrm{im}Q}=_M\left(^0\right)^{}`$. To complete the landscape, we of course define $`_{}^2=\overline{\mathrm{im}Q^{}}`$, so that
(61)
$$\mathrm{ker}d_H\mathrm{\Omega }^2V=\mathrm{ker}Q\overline{\mathrm{im}Q^{}}=_+^2_{}^2.$$
Then in vanishing Webster torsion, one has by (38) that
(62)
$$\begin{array}{cc}\hfill Q(D)& =d_HJ(TJ(d_H_H)^2)=Td_H+(d_H_H)^3\hfill \\ & =\mathrm{\Delta }_HQ,\hfill \end{array}$$
on $`\mathrm{ker}d_H\mathrm{\Omega }^2V`$, where $`\mathrm{\Delta }_H=d_H\delta _H`$ is the contact Laplacian on $`\mathrm{\Omega }^3M`$, conjugate to $`\mathrm{\Delta }_H`$ on functions through $`_M`$. This shows that $`D`$ is conjugate to $`\mathrm{\Delta }_H`$ on $`_M\left(^0\right)^{}`$ by $`Q`$, and that $`D`$ preserves the splitting (61). We therefore recover the decomposition of $`\mathrm{spec}(D)`$ in two families (60), but now entirely seen within $`\mathrm{\Omega }^2V`$ :
(63)
$$\mathrm{spec}^{}(D)=\mathrm{spec}^{}(D_{_{}^2=\overline{\mathrm{im}Q^{}}})\mathrm{spec}^{}(D_{_+^2=\mathrm{ker}Q}).$$
where by (62), $`D`$ is conjugate to $`\mathrm{\Delta }_H`$ on $`_H(^0)^{}`$ by $`Q`$.
The space $`_{}^2`$ is actually a CR invariant, as is $`_+^2`$. Indeed $`\mathrm{\Delta }_H`$ is surjective on $`\mathrm{\Omega }^3M`$ up to โconstantโ $`3`$-forms $`C\theta d\theta `$; as $`Q^{}`$ is zero on these,
$`_{}^2`$ $`=\overline{\mathrm{im}Q^{}}=\overline{\mathrm{im}Q^{}\mathrm{\Delta }_H}`$
$`=\overline{\mathrm{im}DJ\delta _H},\text{by}(\text{62}),`$
$`=\overline{\mathrm{im}DJd_H}.`$
We now have two splittings of $`\mathrm{\Omega }^2V\mathrm{im}D`$ : the spectral one
$$\mathrm{im}D=E^+E^{}$$
in the positive and negative eigenspaces of $`D`$, and the CR invariant one given by
$$\mathrm{im}D=(_+^2\mathrm{im}D)_{}^2.$$
It follows from Prop. 7.6, (60) and (61) that, on a CR-Seifert manifold, the pair $`(E^+,E^{})`$ is in Fredholm position with respect to $`(_+^2,_{}^2)`$. More precisely,
$$_+^2=E^+VH^2(M,)\text{and}E^{}=_{}^2V$$
with the finite dimensional space $`V=_+^2E^{}`$. This enlightens the CR meaning of the spectral asymmetry of $`D`$ we are studying here.
Observe however that if the formal definitions of $`_\pm ^2`$ make sense on any $`3`$-dim CR manifold, their use is highly problematic in general. For instance $`_+^2`$ may be empty if $`M`$ does not bound a Stein manifold, while $`E^+`$ and $`E^{}`$ still exist and keep their nice analytic features by hypoellipticity of $`D`$ on $`\mathrm{im}D`$. The previous Fredholm picture then definitely breaks down. Anyway, from the pseudodifferential viewpoint, the projection on $`E^+`$ is a natural quantization of the real part of the Szegรถ projector on holomorphic $`(2,0)`$-forms, as seen at the Heisenberg symbolic level, see e.g \[5, Chap 4\] for more details on this notion.
We now come back to the comparison between the Riemannian and contact spectra. In (60), we can proceed as in (59) by โfilling the holesโ in the spectrum of $`\mathrm{\Delta }_H`$ on $`^0`$. From (56) we still have $`\mathrm{\Delta }_H=iT`$ on CR functions, and this leads to the following decomposition:
(64)
$$\mathrm{spec}^{}(D)=\mathrm{spec}^{}(\mathrm{\Delta }_H)2\times \mathrm{spec}^{}(iT_{^{2,0}})2\times \mathrm{spec}^{}(iT_{\mathrm{ker}\overline{}_b}).$$
###### 8.1 Remark.
In a slightly more tricky way, one can add $`\mathrm{spec}^{}(\mathrm{\Delta }_H)`$ to both sides of (64): the operator $`\mathrm{\Delta }_H`$ on functions is conjugate to $`\mathrm{\Delta }_H=d_H\delta _H`$ on $`\mathrm{\Omega }^3M`$ and, wedging by $`\theta `$, to $`\delta _Hd_H`$ on $`\mathrm{\Omega }^2V`$. The spectrum of the contact Laplacian
$$\mathrm{\Delta }_2=D+\delta _Hd_H\text{on }\mathrm{\Omega }^2V$$
(see section 9 for more on this one) appears then in a very symmetric manner, namely
(65)
$$\begin{array}{cc}\hfill \mathrm{spec}^{}(\mathrm{\Delta }_2)& =\mathrm{spec}^{}(D)\mathrm{spec}^{}(\mathrm{\Delta }_H)\hfill \\ & =\mathrm{spec}^{}(\mathrm{\Delta }_H)\mathrm{spec}^{}(\mathrm{\Delta }_H)\hfill \\ & 2\times \mathrm{spec}^{}(iT^{2,0})2\times \mathrm{spec}^{}(iT\mathrm{ker}\overline{}_b),.\hfill \end{array}$$
This spectral symmetry can also be seen directly. Equation (38) yields
$$\mathrm{\Delta }_2=T_Hd_H\delta _H+\delta _Hd_H=T_H+P$$
on $`\mathrm{\Omega }^2V=\theta \mathrm{\Omega }^1H`$. As $`[_H,T_H]=0`$ while $`_HP=P_H`$, $`\mathrm{\Delta }_2(_HP)=(_HP)\mathrm{\Delta }_2`$ and $`\mathrm{spec}(\mathrm{\Delta }_2)`$ is symmetric except maybe on $`\mathrm{ker}P`$, where $`\mathrm{\Delta }_2=T_H=TJ`$. It is then easily seen that the kernel splits into
$$(\mathrm{ker}P)^{2,0}=^{2,0}\overline{}_{b}^{}{}_{}{}^{1}(_M\mathrm{ker}\overline{}_b),$$
yielding (65).
###### 8.2 Remark.
Let us mention that this decomposition and the spectral symmetry of $`\mathrm{\Delta }_2`$ also hold on contact manifolds of any dimension, in vanishing Tanaka-Webster torsion, see \[44, Prop. 8\]. This leads to the same kind of formulae as (65), with a โresidual spectrumโ given by sum of $`\eta `$-functions counting holomorphic objects.
### Comparison of contact and Riemannian eta invariants
Comparing the spectrum of $`P_\epsilon `$ given by (59) with that of $`D`$ in (64) yields an immediate relation between their $`\eta `$-functions, up to combinations of $`\zeta `$-functions of positive operators:
###### 8.3 Proposition.
On a CR-Seifert manifold,
(66)
$$\eta (P_\epsilon )\eta (D)=\zeta (\mathrm{\Delta }_H)+\zeta (Q_\epsilon ^+)\zeta (Q_\epsilon ^{}),$$
where $`Q_\epsilon ^\pm ={\displaystyle \frac{1}{2\epsilon }}(\pm 1+\sqrt{1+4\epsilon \mathrm{\Delta }_\epsilon })`$, and $`\mathrm{\Delta }_\epsilon =\mathrm{\Delta }_H\epsilon T^2`$ on functions.
Following Definition 3.2, the renormalized $`\eta `$ invariant $`\eta _0(M,\theta )`$ is the constant term in the development of $`\eta (P_\epsilon )(0)=\eta (M,g_\epsilon )`$ in powers of $`\epsilon `$. It is then immediately extracted from (66) as follows:
(67)
$$\eta _0(M,\theta )=\eta (D)(0)+\zeta (\mathrm{\Delta }_H)(0)+\zeta _0(Q),$$
where $`\zeta _0(Q)`$ is the constant term in the development in powers of $`\epsilon `$
(68)
$$\zeta (Q_\epsilon ^+)(0)\zeta (Q_\epsilon ^{})(0)=\underset{i=2}{\overset{2}{}}\zeta _i(Q)\epsilon ^i,$$
which we already know to exist by (12) and (66), since it is the same as that of $`\eta (P_\epsilon )`$ except for the constant term. Moreover, it turns out that $`\zeta _0(Q)`$ can be evaluated without too much harm on *arbitrary* CR manifolds of dimension $`3`$.
###### 8.4 Lemma.
On any $`3`$-dimensional CR manifold,
$$\zeta (Q_\epsilon ^+)(0)=\zeta (Q_\epsilon ^{})(0),$$
and
$$\zeta _0(Q)=\frac{1}{24\pi ^2}_M|\tau |^2\theta d\theta .$$
where $`\tau =\frac{1}{2}J_TJ`$ is the Tanaka-Webster torsion.
###### Proof.
In view of
$$2\epsilon Q_\epsilon ^\pm =\pm 1+\sqrt{1+4\epsilon \mathrm{\Delta }_\epsilon },$$
we consider for $`\lambda 1`$ the family of operators
$$Q(\lambda )=\lambda +\sqrt{1+4\epsilon \mathrm{\Delta }_\epsilon },$$
where actually
$$\epsilon \mathrm{\Delta }_\epsilon =\epsilon \mathrm{\Delta }_H\epsilon ^2T^2=\mathrm{\Delta }_{g_\epsilon }$$
is the standard Laplacian on functions for the rescaled metric $`g_\epsilon =\epsilon ^2\theta ^2+\epsilon ^1\gamma _H`$ we use here.
Seeleyโs classical results infer that $`Q(\lambda )`$ is a smooth family of positive elliptic pseudo-differential operators of order $`1`$, and that their $`\zeta `$-functions
$$P(\lambda )(s):=\zeta (\lambda +\sqrt{1+4\mathrm{\Delta }_{g_\epsilon }})(s)$$
are meromorphic with possibly simple poles at $`s=1`$, $`2`$ and $`3`$. According to \[4, Prop. 2.9\] or \[24, Lemma 1.10.2\] one can differentiate $`P(\lambda )(s)`$ with respect to $`\lambda `$ to get
$$\frac{d}{d\lambda }P(\lambda )(s)=sP(\lambda )(s+1).$$
Therefore $`{\displaystyle \frac{d^4}{d\lambda ^4}}P(\lambda )(0)=0`$ since $`P(\lambda )`$ is regular at $`s=4`$, and $`P(\lambda )(0)`$ is a polynomial of degree $`3`$ in $`\lambda `$:
(69)
$$P(\lambda )=\zeta ((1+4\mathrm{\Delta }_{g_\epsilon })^{1/2})(0)\lambda R_1+\lambda ^2\frac{R_2}{2}\lambda ^3\frac{R_3}{3},$$
where $`R_0=\zeta (\sqrt{1+4\mathrm{\Delta }_{g_\epsilon }})(0)`$ and $`R_n`$ for $`n>0`$ stands for the residue at $`s=n`$ of
$$\zeta (\sqrt{1+4\mathrm{\Delta }_{g_\epsilon }})(s)=\zeta (1+4\mathrm{\Delta }_{g_\epsilon })(s/2).$$
Actually these residues are related to the development of the heat kernel of $`\mathrm{\Delta }_{g_\epsilon }`$ on functions in a simple way. Let
$$\mathrm{Tr}(\mathrm{e}^{t\mathrm{\Delta }_{g_\epsilon }})\stackrel{t0^+}{}\frac{a_0(g_\epsilon )}{t^{3/2}}+\frac{a_2(g_\epsilon )}{t^{1/2}}+\mathrm{}.$$
According to \[24, Theorem 4.8.18d\], the constants are computed in terms of the volume and the Riemannian scalar curvature of $`g_\epsilon `$ as:
(70)
$$a_0(g_\epsilon )=\frac{\mathrm{Vol}(M,g_\epsilon )}{(4\pi )^{3/2}}\text{and}a_2(g_\epsilon )=\frac{1}{6(4\pi )^{3/2}}_M\mathrm{Scal}(g_\epsilon )d\mathrm{vol}_{g_\epsilon }.$$
This yields
$$\mathrm{Tr}(\mathrm{e}^{t(1+4\mathrm{\Delta }_{g_\epsilon })})=e^t\mathrm{Tr}(\mathrm{e}^{4t\mathrm{\Delta }_{g_\epsilon }})\frac{a_0(g_\epsilon )}{8t^{3/2}}+\frac{4a_2(g_\epsilon )a_0(g_\epsilon )}{8t^{1/2}}+\mathrm{},$$
and by Mellinโs transform \[24, Lemma 1.10.1\],
$$\mathrm{\Gamma }(s/2)\zeta (1+\mathrm{\Delta }_{g_\epsilon })(s/2)=\frac{a_0(g_\epsilon )}{4(s3)}+\frac{4a_2(g_\epsilon )a_0(g_\epsilon )}{4(s1)}+h(s),$$
with $`h`$ holomorphic for $`\mathrm{Re}s>1`$. Hence
$$\zeta ((1+4\mathrm{\Delta }_{g_\epsilon })^{1/2})(0)=0$$
as this is the only way to cancel the simple pole of the Gamma function at $`s=0`$, and
$$R_2=0,$$
(because the Gamma function does not vanish at $`s=2`$ and the r.h.s. has no pole at this point) so that $`P(\lambda )`$ is an odd polynomial. This gives $`P(1)=P(1)`$ or, equivalently,
$$\zeta (Q_\epsilon ^+)(0)=\zeta (Q_\epsilon ^{})(0)$$
as announced. Moreover one has
$$R_1=\frac{4a_2(g_\epsilon )a_0(g_\epsilon )}{4\sqrt{\pi }}\mathrm{and}R_3=\frac{a_0(g_\epsilon )}{2\sqrt{\pi }},$$
and thus by (69) and (70)
(71)
$$\begin{array}{cc}\hfill \zeta (Q_\epsilon ^+)(0)& =R_1R_3/3\hfill \\ & =\frac{1}{\sqrt{\pi }}(\frac{a_0(g_\epsilon )}{12}a_2(g_\epsilon ))\hfill \\ & =\frac{1}{48\pi ^2\epsilon ^2}\left(\frac{1}{2}_M\theta d\theta _M\mathrm{Scal}(g_\epsilon )\theta d\theta \right).\hfill \end{array}$$
The Riemannian curvature of $`g_\epsilon `$ can be developed in powers of $`\epsilon `$ using the links between Tanaka-Webster and Levi-Civita connections underlined in (13). According to *e.g.* \[44, p 318\], one finds in dimension $`3`$ that
$$\mathrm{Scal}(g_\epsilon )=\frac{1}{2}+2\epsilon R\epsilon ^2|\tau |^2,$$
where $`R`$ and $`\tau `$ are Tanaka-Webster curvature and torsion. The constant term in the full development of $`\zeta (Q_\epsilon ^+)`$ is then necessarily equal to the integral of $`\frac{1}{48\pi ^2}|\tau |^2`$ on $`M`$. โ
###### 8.5 Remark.
According to (59), $`Q_\epsilon ^+`$ describes the non collapsing spectrum of $`d_\epsilon `$, on Seifert-CR manifolds. We have seen that this spectrum only contributes by a local expression $`\zeta (Q_\epsilon ^+)(0)`$ to $`\eta (d_\epsilon )`$. We expect this to hold in the general case. Indeed on any CR manifold, the non-collapsing spectrum is always strictly positive, since it converges to $`1`$ and $`d_\epsilon `$ has no spectral flow. It therefore always contributes through a *zeta* function, whose value at $`0`$ is local for a wide class of operators.
### A computation of $`\eta _0`$.
The previous Lemma 8.4, together with the spectral decomposition (59), leads to a general computation of the renormalized $`\eta `$-invariant on all CR-Seifert manifolds, including the still missing exceptional cases of section 5. Indeed, one has
$$\zeta ^{}(Q_\epsilon ^+)\zeta ^{}(Q_\epsilon ^{})=\zeta (Q_\epsilon ^+)\zeta (Q_\epsilon ^{})+1,$$
since $`0`$ belongs to $`\mathrm{spec}(Q_\epsilon ^{})`$ with multiplicity $`1`$ (corresponding to the constant functions). It follows then from (59) that
(72)
$$\eta _0(d)=\eta _{\mathrm{ad}}(d)=1+2(\eta (iT_{^{2,0}})(0)\eta (iT_{\mathrm{ker}\overline{}_b})(0)).$$
These holomorphic counting functions can be nicely expressed as dimensions of spaces of sections on adequate orbifold line bundles over the basis orbifold Riemann surface, which in turn are easily computed with the help of Riemann-Roch-Kawasakiโs theorem . Note that this has already been observed in the adiabatic setting and constant curvature by L. Nicolaescu in \[39, Sec. 1\]. We give below only a short description of the computation, and refer to for more details.
Following section 5, the CR-Seifert manifold $`M`$ may be seen as the unit circle bundle of some orbifold line bundle $`L`$ over $`\mathrm{\Sigma }`$, with singular data $`(\alpha _i,\rho _i,\beta _i)`$ at points $`m_i\mathrm{\Sigma }`$. Let $`K_\mathrm{\Sigma }=\mathrm{\Lambda }^{1,0}T^{}\mathrm{\Sigma }`$ denotes the orbifold canonical bundle of $`\mathrm{\Sigma }`$. Now, given a Fourier component $`iT=n`$, the space of CR functions $`f`$ such that $`f\phi _t=e^{int}f`$ are interpreted as the space of holomorphic sections of $`L^n`$, and we denote by $`h_0(L^n)`$ its dimension. Moreover the space of holomorphic forms $`\sigma `$ in the canonical CR bundle $`K_M\theta K_\mathrm{\Sigma }L`$ such that $`iT\sigma =n\sigma `$ may be seen as the space of holomorphic sections of $`K_\mathrm{\Sigma }L^n`$, *i.e.* $`(1,0)`$-holomorphic forms in $`L^n`$. Let $`h_1(L^n)`$ denotes its dimension. Hence we get
(73) $`\eta (iT_{^{2,0}})(s)\eta (iT_{\mathrm{ker}\overline{}_b})(s)`$ $`={\displaystyle \underset{n^{}}{}}\mathrm{sgn}(n){\displaystyle \frac{h_0(L^n)h_1(L^n)}{|n|^s}}`$
$`={\displaystyle \underset{n^{}}{}}\mathrm{sgn}(n){\displaystyle \frac{\chi _\overline{}(L^n)}{|n|^s}}.`$
Following the method in \[39, Sec. 1\], this sum can be computed explicitly using Riemann-Roch-Kawasaki theorem (extension of the classical Riemann-Roch to the orbifold case) . Using the (rational) orbifold Euler characteristic $`\chi `$ of the base $`\mathrm{\Sigma }`$ and the (rational) degree $`d`$ of $`L`$, it reads
(74)
$$\chi _\overline{}(L^n)=\frac{\chi }{2}nd+\underset{i}{}\frac{1}{2}\left(1\frac{1}{\alpha _i}\right)\left\{\frac{n\beta _i\rho _i^{}}{\alpha _i}\right\},$$
where $`\{x\}=x[x]`$ denotes the fractional part of $`x`$, and $`\rho _i^{}`$ is the inverse of $`\rho _i`$ mod. $`\alpha _i`$. This purely topological formula holds true, irrespective of the curvature value. The result should then be the same in the constant and non-constant curvature cases, so that Ouyangโs formula (26) for $`\eta _0`$ holds true on any CR-Seifert manifold.
To get explicitly the formula, one can argue as follows: the constant terms in (74) do not contribute to the sum (73), whereas
$$\underset{n^{}}{}d|n|^{s+1}=2d\zeta (s1)$$
has value $`\frac{d}{6}`$ at $`s=0`$. The Dedekind-Rademacher sums $`s(\alpha _i,1,\beta _i\rho _i^{})=s(\alpha _i,\rho _i,\beta _i)`$ appear from the periodic orbifold contribution in (74), as in Nicolaescuโs work using \[39, Proposition 1.4\]. Inserting in (72) leads to the desired expression.
###### 8.6 Remark.
This last computation shows that Theorem 1.4 could have been proved in a quicker way on constant curvature CR-Seifert manifolds: applying the previous formulae and using the computation of $`\zeta (\mathrm{\Delta }_H)(0)`$ given below leads to an expression for $`\eta (D)`$ that can be compared directly to Ouyangโs formula for $`\eta _0`$. We have however omitted this proof since the links between $`\eta (D)`$ and $`\eta _0`$ proved in this way would have appeared as the result of a possibly completely fortuitous or miraculous equality between explicitly known numerical expressions. On the contrary, our proof stresses the fact that the relation between $`D`$ and $`d`$ is deeply rooted in the nature of CR geometry and the diabatic limit. Moreover, it applies to *the whole family* of CR-Seifert manifolds, irrespective of their curvature, and especially the exceptional cases that do not admit constant curvature contact forms.
We now complete the comparison between $`\eta _0`$ and the contact $`\eta `$-invariant $`\eta (D)`$.
###### 8.7 Theorem.
Let $`M`$ be a CR-Seifert manifold. Then,
(75)
$$\eta _0(M,\theta )=\eta (D)(0)+\zeta (\mathrm{\Delta }_H)(0)$$
with
(76)
$$\zeta (\mathrm{\Delta }_H)(0)=\frac{1}{512}_MR^2\theta d\theta .$$
###### Proof.
From Proposition 8.3 and Lemma 8.4 it remains to compute $`\zeta (\mathrm{\Delta }_H)(0)`$. The development of the heat kernel $`\mathrm{e}^{t\mathrm{\Delta }_H}`$ of the Kohn Laplacian $`\mathrm{\Delta }_H`$ has been studied by Beals, Greiner and Stanton in \[6, Theorem 7.30\]. On any CR manifold of dimension $`3`$,
$$\mathrm{Tr}(\mathrm{e}^{t\mathrm{\Delta }_H})\underset{n=0}{\overset{\mathrm{}}{}}t^{n2}b_n(M,\theta )\text{as}t0^+,$$
where $`b_n(M,\theta )`$ are integrals over $`M`$ of polynomials of covariant derivatives of Tanaka-Webster curvature and torsion. Mellinโs transform yields again
$$\mathrm{\Gamma }(s)\zeta (\mathrm{\Delta }_H)(s)=\underset{nN}{}\frac{b_n(M,\theta )}{s2+n}+h_N(s)$$
with $`h_N`$ holomorphic for $`\mathrm{Re}s>N2`$, and hence
$$\zeta (\mathrm{\Delta }_H)(0)=b_2(M,\theta ).$$
As $`\zeta (\mathrm{\Delta }_H)(0)`$ stays unchanged when $`\theta `$ becomes $`k\theta `$, one must have $`b_2(M,k\theta )=b_2(M,\theta )`$, and the same argument as in Lemma 4.1 gives that
$$b_2(M,\theta )=C_1_MR^2\theta d\theta +C_2_M|\tau |^2\theta d\theta ,$$
for some constants $`C_1`$, $`C_2`$.
Thanks to N. Stantonโs work it is possible to determine $`C_1`$ on the sphere $`๐^3`$. Indeed, let $`L=4\mathrm{\Delta }_H+R`$ be the CR-conformal Laplacian on $`๐^3`$. Stanton states in \[49, Theorem 4.34\] that for the contact form $`\theta =i\overline{}r=\frac{i}{2}(z^1d\overline{z}^1+z^2d\overline{z}^2)`$
$$\mathrm{Tr}(\mathrm{e}^{tL})=\frac{\pi ^2}{256t^2}+O\left(\frac{1}{t^2}\mathrm{e}^{\pi ^2/4t}\right)\text{as}t0^+.$$
Now Tanaka-Webster curvature $`R=4`$ here, so that the heat development of $`\mathrm{\Delta }_H`$ is
$`\mathrm{Tr}(\mathrm{e}^{t\mathrm{\Delta }_H})=\mathrm{e}^t\mathrm{Tr}(\mathrm{e}^{tL/4})=\mathrm{e}^t{\displaystyle \frac{\pi ^2}{16t^2}}+O\left({\displaystyle \frac{1}{t^2}}\mathrm{e}^{\pi ^2/4t}\right),`$
and the constant term is $`b_2(M,\theta )=\frac{\pi ^2}{32}`$. Hence
$$\zeta (\mathrm{\Delta }_H)(0)=\frac{\pi ^2}{32}=C_1_{๐^3}R^2\theta d\theta =16\pi ^2C_1$$
yields $`C_1=\frac{1}{32\times 16}`$ on the sphere, hence on any CR-Seifert manifold. โ
Putting together this last result and Theorem 1.1 leads to Corollary 1.5.
## 9. The contact and the modified contact $`\eta `$-invariants
We first begin by showing existence of the contact $`\eta `$-invariant in dimension $`3`$. It follows mostly the classical method of Chapter 1 of , using pseudo-differential calculi developed on contact manifolds. As a consequence, we shall put below the emphasis mainly on the steps where the hypoelliptic context introduces differences with the well-known elliptic theory.
###### 9.1 Theorem.
Let $`(M,H,J)`$ be a compact $`3`$-dimensional strictly pseudoconvex CR manifold endowed with a compatible contact form $`\theta `$ and the associated metric $`g_1=\theta ^2+d\theta (,J)`$. Then the series
$$\eta (D)(s)=\mathrm{Tr}^{}(D|D|^{(s+1)})=\underset{\lambda _i\mathrm{spec}(D)\{0\}}{}\frac{\lambda _i}{|\lambda _i|^{s+1}}$$
converges absolutely for $`\mathrm{Re}s>2`$, and has an meromorphic extension with possible simple poles at $`s=2n/2`$ for $`n`$. Moreover $`\eta (D)(s)`$ is regular at $`s=0`$; its value $`\eta (D)(0)`$ is the *contact $`\eta `$-invariant*.
###### Proof.
From the two Laplacians
$$\mathrm{\Delta }_2=D+\delta _Hd_H\text{on}\mathrm{\Omega }^2V\text{and}\mathrm{\Delta }_3=d_H\delta _H\text{on}\mathrm{\Omega }^3M,$$
are maximally hypoelliptic (be careful: $`\mathrm{\Delta }_3`$ is nonnegative, but $`\mathrm{\Delta }_2`$ is not, despite the notation). This means that they control two horizontal derivatives in $`L^2`$ norms (and one vertical derivative). By the associated Sobolev embeddings, their resolvents are compact and their spectra are discrete. By orthogonality and conjugation, the non-zero spectrum of $`\mathrm{\Delta }_2`$ splits into
(77)
$$\mathrm{spec}^{}(\mathrm{\Delta }_2)=\mathrm{spec}^{}(D)\mathrm{spec}^{}(\mathrm{\Delta }_3),$$
and $`D`$ has discrete pure point spectrum with finite multiplicities on $`\mathrm{im}D`$. Sobolev embeddings also yields that $`(i+\mathrm{\Delta }_2)^n`$, $`(i+\mathrm{\Delta }_3)^n`$ are trace class for $`n`$ large enough, hence the same for $`(D)^n`$. The series $`\eta (D)(s)`$ is then well defined and holomorphic for $`\mathrm{Re}s`$ large.
Getting more information on $`\eta `$ relies in the Riemannian (elliptic) case on the use of the classical pseudo-differential calculus for elliptic operators. Such a symbolic calculus has also been developed on contact manifold by Beals, Greiner and Stanton in or Taylor in , a concise account may also be found in . The symbols of the hypoelliptic operators $`\mathrm{\Delta }_2`$ and $`\mathrm{\Delta }_3`$ are invertible in this calculus: this follows from \[29, Lemmas 5.18, 5.19\], or else by observing that in dimension $`3`$ their principal symbols are sums of invertible Folland-Stein ones.
The parameter calculus adapted to the Heisenberg setting developed in propositions 5.20 to 5.26 of yields pseudo-differential approximations $`R(\lambda )`$ of the resolvents $`((\mathrm{\Delta }_2)^2\lambda )^1`$, when $`\lambda ^+`$. This uses the classical iteration process described in \[24, p. 51\] or \[48, Sec. 9.1\] for instance, where the standard pseudo-differential symbolic product has to be replaced by the Heisenberg one, see . The symbol of these $`R(\lambda )`$ are universal expressions involving the symbol of $`(\mathrm{\Delta }_2)^2\lambda `$, its inverse, and tensorial expressions of the Webster-Tanaka curvature and its derivatives.
Then, as explained in \[24, Sec. 1.7\], $`R(\lambda )`$ can be used in place of $`((\mathrm{\Delta }_2)^2\lambda )^1`$ in the contour integral
$$\mathrm{\Delta }_2\mathrm{e}^{t(\mathrm{\Delta }_2)^2}=\frac{1}{2i\pi }_\gamma \mathrm{e}^{t\lambda }\mathrm{\Delta }_2(\mathrm{\Delta }_2^2\lambda )^1๐\lambda ,$$
with $`\gamma ^+`$ the correctly oriented boundary of the cone $`\{\mathrm{Im}\lambda \mathrm{Re}\lambda +1\}`$, in order to get good approximations of $`\mathrm{\Delta }_2\mathrm{e}^{t(\mathrm{\Delta }_2)^2}`$ when $`t`$ goes to $`0`$. Following Lemma 1.7.7 of , homogeneity arguments then easily lead to the asymptotic development of $`\mathrm{Tr}(\mathrm{\Delta }_2\mathrm{e}^{t\mathrm{\Delta }_2^2})`$ when $`t0^+`$. Namely,
(78)
$$\mathrm{Tr}(\mathrm{\Delta }_2\mathrm{e}^{t\mathrm{\Delta }_2^2})\underset{n=0}{\overset{\mathrm{}}{}}t^{(n6)/4}R_n(M,\theta ),$$
where $`R_n(M,\theta )`$ are integrals over $`M`$ of universal polynomials in Tanaka-Webster curvature and covariant derivatives (with respect to the classical elliptic development given in \[24, Lem 1.7.7\], the only changes here concern the powers of $`t`$: this is due to the fact that, in the Heisenberg calculus, horizontal directions have weight $`1`$, while $`T`$ is of weight $`2`$. For instance, this implies that the โHeisenberg-dimensionโ of $`M`$ is $`4`$ instead of $`3`$).
###### 9.2 Remark.
Another more direct track, if steeper, also leads to such kernel developments. One can follow Beals-Greiner-Stantonโs approach to heat kernels asymptotics in the contact setting. In they have extended their symbolic calculus on $`M\times `$ to include the heat operator $`_t+P`$ for some positive sub-Laplacians $`P`$. They show that in the case $`P`$ is a positive Folland-Stein type operator, one can inverse the symbol of $`_t+P`$ inside this calculus, which gives rather directly developments like (78) for $`\mathrm{Tr}(Q\mathrm{e}^{tP})`$ from the symbol of $`Q(_t+P)^1`$, see also \[23, Sec 4\]. By R. Pongeโs recent work , this approach leads to a relatively simple proof of the index theorem, and also applies to more general positive hypoelliptic $`P`$ as $`(\mathrm{\Delta }_2)^2`$.
Let us now complete the proof of Theorem 9.1. Mellin transform and the functional calculus relate the asymptotic development in small time of the heat kernel to $`\eta `$ and $`\zeta `$ functions \[24, Section 1.10\]. In particular, \[24, p 81\] and (78) yield:
$$\eta (\mathrm{\Delta }_2,s)\mathrm{\Gamma }((s+1)/2)=\underset{n=0}{\overset{N}{}}\frac{4}{2s+n4}R_n(M,\theta )+h_N(s)$$
where $`h_N`$ is an holomorphic function for $`s>2N/2`$. Hence we get the required meromorphic extension of $`\eta (\mathrm{\Delta }_2)(s)`$. The same technique applies to $`\mathrm{\Delta }_3`$ on $`\mathrm{\Omega }^3M`$, but this is a positive operator whose heat kernel development has been extensively treated in \[6, Theorem 7.30\]: the $`\eta `$-function is here a $`\zeta `$-function which is regular at $`s=0`$.
Using the spectral decomposition (77), we get that $`\eta (D)(s)`$ is meromorphic with $`s=0`$ being possibly a simple pole. It remains to show that this function is regular at $`s=0`$. We first note that the value of the residue of $`\eta (D)`$ at $`s=0`$ is $`2R_4(M,\theta )`$. It is easily seen in (38) that $`D`$ becomes $`kD`$ in the contact rescaling $`\theta k\theta `$. Therefore, $`\eta (D_{k\theta })(s)=k^s\eta (D_\theta )(s)`$ and
$$R_4(M,k\theta )=R_4(M,\theta ).$$
Following the proof of Lemma 4.1, this implies that, *in dimension $`3`$*,
(79)
$$R_4(M,\theta )=C_1_MR^2\theta d\theta +C_2_M|\tau |^2\theta d\theta $$
where $`R`$ and $`\tau `$ are Tanaka-Webster curvature and torsion and $`C_1`$, $`C_2`$ are universal constants.
The residue is moreover invariant under smooth deformation of the pseudohermitian and CR structures (*i.e.* both $`\theta `$ and $`J`$): as underlined in \[24, Lemma 1.10.2\] this general feature stems from the existence of a local variation formula for $`\eta `$-functions, namely in the absence of spectral flow here:
$$\dot{\eta }(\mathrm{\Delta }_2)(s)=s\mathrm{Tr}(\dot{\mathrm{\Delta }}_2\mathrm{\Delta }_2^{(s+1)/2}).$$
The point here is that the trace on the right has a meromorphic extension coming from the development of $`\mathrm{Tr}(\dot{\mathrm{\Delta }}_2\mathrm{e}^{t(\mathrm{\Delta }_2)^2})`$, but the possible simple pole at $`s=0`$ is actually cancelled out by the $`s`$ in front of the whole expression.
The conclusion is that the integrals in (79) have to be independent of variations of $`\theta `$ and $`J`$, and this implies $`C_1=C_2=0`$: indeed, the variations of $`R^2`$ and $`|\tau |^2`$ when $`\theta \theta _f=\mathrm{e}^{2f}\theta `$ have been computed in \[34, Sec. 5\]. One finds that
(80)
$$\frac{d}{df}(R^2\theta d\theta )=8R(\mathrm{\Delta }_Hf)\theta d\theta $$
while (if $`\tau =A_{11}\theta ^1\theta ^1`$)
(81)
$$\frac{d}{df}(|\tau |^2\theta d\theta )=2i(A_{\overline{1}\overline{1}}f_{,11}A_{11}f_{,\overline{1}\overline{1}})\theta d\theta .$$
After integration by parts, this yields
(82)
$$\frac{d}{df}R_4(M,\theta )=8C_1_Mf\mathrm{\Delta }_HR\theta d\theta +\mathrm{\hspace{0.17em}2}iC_2_Mf(A_{\overline{1}\overline{1},11}A_{11,\overline{1}\overline{1}})\theta d\theta .$$
Testing on a circle bundle (with vanishing torsion) over a Riemann surface of non constant curvature cancels out $`C_1`$. General expression for torsion of hypersurfaces in \[52, Sec. 4\] shows that $`A_{\overline{1}\overline{1},11}A_{11,\overline{1}\overline{1}}`$ does not vanish identically: actually, following the only Bianchi identity of order $`2`$ between $`R`$ and $`\tau `$ in dimension $`3`$ is $`R_{,0}=A_{11,\overline{1}\overline{1}}+A_{\overline{1}\overline{1},11}`$, which does not occur in (82) so that $`C_2=0`$. โ
###### 9.3 Remark.
The contact-de Rham complex exists on contact manifolds of any dimension, and the contact-signature operator $`D`$ is still self-adjoint in dimension $`4n1`$. Therefore the properties of $`\eta (D)(s)`$ stated in Theorem 9.1 make sense on contact manifolds of any dimension. Most of the previous discussion, and its conclusions, still applies, but the last argument about the regularity at $`s=0`$ of $`\eta (D)`$. The residue is still both a contact invariant, independent of the choices of $`\theta `$ and $`J`$, and an integral of some universal pseudohermitian polynomial of the right weight. But many possibilities are now left, which cannot be so easily analysed (even in the next relevant dimension $`7`$, the algebra becomes quite complicated). At the present time, one still ignores whether this residue always vanishes or not.
### The CR invariant correction of $`\eta (D)`$
Having now a well-defined object at hand, we can proceed to the construction of a *modified contact $`\eta `$-invariant*.
###### 9.4 Theorem.
There exists a unique choice of universal constants $`C_1`$ and $`C_2`$ such that, for any compact strictly pseudoconvex CR $`3`$-manifold $`M`$, the following pseudohermitian invariant
(83)
$$\overline{\eta }(D)=\eta (D)+C_1_MR^2\theta d\theta +C_2_M|\tau |^2\theta d\theta ,$$
formed from a contact form $`\theta `$, its Tanaka-Webster curvature $`R`$ and torsion $`\tau `$, is in fact a *CR invariant* of $`M`$, which we shall call the *modified contact* $`\eta `$-*invariant*.
The key point for the proof of Theorem 9.4 is the following: on an oriented CR $`3`$-manifold $`M`$, the space of adapted contact forms for a given CR structure (let us denote it by $`\mathrm{\Theta }`$) is contractible and non-empty. Then, for a CR invariant, being CR invariant simply means being independent of the choice of the contact form, *i.e.* having a vanishing derivative in the direction of any variation in $`\theta `$.
Using the analysis above, we get that $`\eta (D)`$, seen as a function on the space $`\mathrm{\Theta }`$ of contact forms adapted to a given CR structure, has the following features :
1. $`\eta (D_{k\theta })=\eta (D_\theta )`$ for any positive $`k`$;
2. its derivative is local: if $`\theta _t=(1+tf)\theta `$ is a small variation of contact forms,
$$\frac{d}{dt}\eta (D_{\theta _t}){}_{t=0}{}^{}=_Mf_\theta \theta d\theta ,$$
where $`_\theta `$ is a local pseudohermitian invariant of $`\theta `$ built algebraically and *universally* from a finite jet of $`\theta `$ and its Tanaka-Webster curvature $`R`$ and torsion $`\tau `$.
One then deduces from (i) and (ii) that, necessarily,
(84)
$$_{k\theta }=k^4_\theta ,$$
and moreover
(85)
$$_M_\theta \theta d\theta =0.$$
Said otherwise, $`_\theta `$ is of weight $`4`$ and vanishing integral. One can then remark a basic fact:
###### 9.5 Lemma.
Let $`\alpha `$ be a smooth *closed*, and *real* $`1`$-form on $`\mathrm{\Theta }`$ where $`T_\theta \mathrm{\Theta }`$ is identified to the space of functions on $`M`$ through $`f\frac{d}{dt}(1+tf)\theta `$. If $`\alpha `$ is of the type
(86)
$$\alpha _\theta :fC^{\mathrm{}}(M)\alpha _\theta (f)=_Mf๐_\theta \theta d\theta $$
where $`๐_\theta `$ is an universal local pseudohermitian invariant of a finite jet of $`\theta `$ of weight $`4`$ and vanishing integral, then $`\alpha `$ is a linear combination of the derivatives in $`\theta `$ of
$$_MR^2\theta d\theta \text{and}_M|\tau |^2\theta d\theta .$$
###### Proof.
We argue as in section 4, classifying local pseudo-hermitian invariants that are real and of weight $`4`$. We have seen that the sole possibilities are:
$$\begin{array}{cc}\hfill R^2,|\tau |^2,& R_{,0}=A_{11,\overline{1}\overline{1}}+A_{\overline{1}\overline{1},11}\text{(Bianchi identity)},\hfill \\ & \mathrm{\Delta }_HR,i(A_{11,\overline{1}\overline{1}}A_{\overline{1}\overline{1},11}).\hfill \end{array}$$
The first two expressions have non-vanishing integrals in general, they then have to be forgotten. From (80) the fourth is the variation of $`\frac{1}{8}_MR^2\theta d\theta `$, whereas from (81) the fifth is the variation of $`\frac{1}{2}_M|\tau |^2\theta d\theta `$.
We check that the third one does not yield a closed form. According to \[34, Sec. 5\], a change of contact form $`\theta \theta _f=e^f\theta `$ induces the following changes
$$R_f=e^f(R+2\mathrm{\Delta }_Hf2|f_{,\overline{1}}|^2)\mathrm{and}T_f=e^f(T+if_1Z_{\overline{1}}if_{\overline{1}}Z_1),$$
and therefore
$$\frac{d}{df}(R_{,0}\theta d\theta )=\left(f_{,0}R+if_{,1}R_{,\overline{1}}if_{,\overline{1}}R_{,1}+2(\mathrm{\Delta }_Hf)_{,0}\right)\theta d\theta .$$
When restricted on the sphere $`๐^3`$ with its constant curvature pseudohermitian structure this gives
$`{\displaystyle _M}(g{\displaystyle \frac{d}{df}}f{\displaystyle \frac{d}{dg}})(R_{,0}\theta d\theta )`$ $`=2{\displaystyle _M}((\mathrm{\Delta }_Hf)_{,0}g(\mathrm{\Delta }_Hg)_{,0}f)\theta d\theta `$
$`=4{\displaystyle _M}(\mathrm{\Delta }_Hf)(T.g)\theta d\theta .`$
This expression does not vanish identically: for instance when taking any non $`T`$-invariant function $`g`$ and $`f`$ such that $`\mathrm{\Delta }_Hf=T.g`$. This completes the proof. โ
This shows Theorem 9.4, exhibiting a new CR invariant
(87)
$$\overline{\eta }(D)=\eta (D)+C_1_MR^2\theta d\theta +C_2_M|\tau |^2\theta d\theta .$$
Uniqueness in the choice of the constants is obtained because no linear combination in the integrals of $`R^2`$ and $`|\tau |^2`$ can be a CR invariant. โ
###### 9.6 Remark.
An analogous line of reasoning yields: there exists a universal constant $`C^{}`$ such that, for any compact strictly pseudoconvex Cauchy-Riemann $`3`$-manifold $`M`$,
(88)
$$\overline{\eta }(D)C^{}\nu (M)$$
is a *contact* invariant, *i.e.* is independent of the choice of the complex structure. The proof (left to the reader) consists in proving that the only tensorial choice for the differential of $`\overline{\eta }`$ is (up to some multiplicative constant) the Cartan curvature like in (28) and (29).
Of course, in view of the relation (7) between $`\nu `$ and $`\eta (D)`$ in the CR-Seifert case, one expects that the constants $`C^{}`$ above and and $`C_1`$ in Theorem 9.4 should be respectively $`\frac{1}{3}`$ and $`(\frac{1}{512}\frac{1}{48\pi ^2})`$, but the case of CR-Seifert manifolds is not sufficient to determine them. The best one can get is the following: it has already been remarked earlier that the value of the renormalized $`\eta `$-invariant $`\eta _0`$ is purely topological on CR-Seifert manifolds. Keeping the contact form fixed, this means that it has to be independent of the complex structure. As $`\eta (D)=\eta _0\frac{1}{512}R^2\theta d\theta `$ and
$$\overline{\eta }C^{}\nu =(1+3C^{})\eta _0+(C_1\frac{1}{512}\frac{C^{}}{16\pi ^2})R^2\theta d\theta $$
must be a contact invariant, this implies that
$$C_1\frac{1}{512}\frac{C^{}}{16\pi ^2}=0,$$
since the integral of $`R^2`$ has non-zero variations with respect to the complex structures.
Guessing the values of $`C`$ in Conjecture 1.6 and $`C_2`$ in Theorem 9.4 seems much harder. Having a precise value for them would (for instance) involve a precise computation of the spectrum of $`\eta (D)`$ in a case where the torsion does not vanish. This seems difficult to achieve either with our methods, which rely on Fourier decomposition under the circle action, or with classical tools of representation theory, which require a high degree of homogeneity.
Of course, one knows that the derivative of $`\eta (D)`$ is given by algebraic expressions of the jet of the hypoelliptic symbols of the involved operators. However these expressions are so intricate that the constants are only computable this way โin theoryโ, and not in practice.
###### 9.7 Remark.
The same arguments also apply to the renormalized $`\eta `$-invariant $`\eta _0`$ introduced in section 3, instead of $`\eta (D)`$. This explains *a priori* the existence of some local correction of $`\eta _0`$ leading to a CR invariant, itself related (up to some contact invariant) to a multiple of $`\nu `$; this might be compared with Lemma 4.1.
## 10. Proof of the corollaries
Corollaries 1.7 and 1.9 rely on the formula discovered by the first and second authors \[11, Theorem 1.2\]: for any Einstein asymptotically hyperbolic manifold $`(N^4,g)`$,
(89)
$$\frac{1}{8\pi ^2}_N\left(3|W^{}|^2|W^+|^2+\frac{1}{24}\mathrm{Scal}^2\right)\chi (N)+3\tau (N)=\nu (M).$$
For complex hyperbolic surfaces, the integral term is zero. If $`\overline{N}`$ is smooth, with $`M`$ as the only end, then the topological contributions always are integers. Corollary 1.7 is then proved.
It is instructive to check the results for a holomorphic disk bundle over a hyperbolic Riemann surface $`\mathrm{\Sigma }`$, with $`M`$ as its boundary. Clearly one has $`\chi (N)=\chi (\mathrm{\Sigma })=\chi `$ and $`\tau (N)=1`$. If $`N`$ carries a complex hyperbolic metric with $`M`$ as its boundary at infinity, then corollary 1.7 gives the equation
$$\chi \mathrm{\hspace{0.17em}3}\tau =\nu (M)=d+3+\frac{\chi ^2}{4d}$$
and the only solution is $`d=\frac{\chi }{2}`$. We then recover the well-known fact that the only disk bundles carrying a complex hyperbolic metric are the square roots of the (complex) tangent bundle.
Corollary 1.9 is again a direct consequence of (89), since for a Kรคhler-Einstein metric, the integral term is non negative. For an Einstein metric, the story is more complicated, but positivity is achieved if solutions of the Seiberg-Witten equations exist, and it is proved in \[43, corollary 31\] that it is a consequence of the nonvanishing of the Kronheimer-Mrowka invariants .
From \[17, Theorem 5.12\], one knows that pseudoconvex complex hyperbolic surfaces $`N`$ have vanishing third homology group $`H_3(N,)`$. Hence no multiple ends can occur, but one expects orbifold singularities or cusps to appear in the interior of a complex hyperbolic filling. The complex hyperbolic cusps can be compactified to yield a complex orbifold surface that we note again $`N`$, by adding at the infinity of each cusp a quotient $`\mathrm{\Sigma }_i`$ of a 2-torus. The Corollaries 1.7 and 1.9 remain true in this case, with the Euler characteristic and the signature of $`N`$ being replaced by their orbifold versions: In case $`\mathrm{}`$ cusps are present, there is an additional contribution in the signature coming from the self-intersection of each 2-torus at infinity. Namely, one has to consider the modified signature \[9, proposition 3.4\]
$$\tau _{\mathrm{cusp}}(N)=\tau (N)\frac{1}{3}\underset{1}{\overset{\mathrm{}}{}}[\mathrm{\Sigma }_i][\mathrm{\Sigma }_i].$$
Of course, Corollary 1.8 is no more true, since the characteristic numbers are now rational; the denominator of $`\nu `$ only gives an hint on the order of the singularities needed to fill $`M`$.
### Explicitation for lens spaces
We now specialize the formula obtained in Corollary 1.3 to the lens space $`L(p,q)`$ obtained as a quotient of the 3-sphere $`๐^3`$ in $`^2`$ by $`/p`$, with its generator acting on $`^2`$ by $`(\mathrm{e}^{\frac{2i\pi }{p}},\mathrm{e}^{\frac{2iq\pi }{p}})`$, where $`q`$ is prime with $`p`$. They are interesting in connection with filling by Einstein metrics, since some of them appear as boundary at infinity of selfdual Einstein metrics . On the other hand, it has been shown that large families of them admit symplectic fillings , so that Corollary 1.9 may be applied to these.
###### 10.1 Proposition.
One has: $`\nu (L(p,q))=\frac{1}{p}+12s(p,q,1)`$.
For sake of comparison, we recall to the interested reader the value of the classical $`\eta `$-invariant on lens spaces with the standard round metric, as computed by Atiyah-Patodi-Singer \[3, Proposition 2.12\]:
(90)
$$\eta (L(p,q))=4s(p,q,1).$$
###### Proof.
For simplicity, we shall assume that $`(q1)`$ is prime with $`p`$ (as a matter of fact this implies that we take $`q1`$), and we leave the general case to the reader. Let us see the 3-sphere as the bundle $`๐ช(1)`$ over the projective line $`P^1`$. The induced action on $`P^1`$ has two fixed points: the two antipodal points, with action of $`/p`$ generated by $`\mathrm{e}^{\pm i2\pi \frac{q1}{p}}`$, and action in the fiber by $`\mathrm{e}^{i\frac{2\pi }{p}}`$ and $`\mathrm{e}^{i2\pi \frac{q}{p}}`$ respectively. Therefore $`L(p,q)`$ is a $`๐^1`$-orbifold bundle over an orbifold projective line with two orbifold points with angle $`\frac{2\pi }{p}`$. The Euler characteristic is $`\chi =\frac{2}{p}`$ and the degree (first Chern number) is $`d=\frac{1}{p}`$. Now Corollary 1.3 and Ouyangโs Theorem 5.2 give the formulae
$`\nu (L(p,q))`$ $`=3+{\displaystyle \frac{2}{p}}12\left(s(p,q1,1)+s(p,1q,q)\right),`$
$`\eta (L(p,q))`$ $`=1{\displaystyle \frac{1}{p}}+4\left(s(p,q1,1)+s(p,1q,q)\right),`$
(note that the extra parameter $`\rho `$ in Theorem 5.2 appears naturally on lens spaces), so that $`\nu (L(p,q))=\frac{1}{p}3\eta (L(p,q))`$. The proposition then follows from (90). โ
### Comparison with the Burns-Epstein invariant
Another interesting point is to compare these results with those obtained by use of the Burns-Epstein $`\mu `$-invariant (it is already suggested at the end of that obstructions follow from computations of $`\mu `$). The $`\mu `$-invariant is defined on strictly pseudoconvex CR 3-manifolds with trivial tangent holomorphic bundle only. Roughly speaking, it comes from Chern-Simons-type constructions (integration of a local formula), whereas the $`\nu `$-invariant is extracted from the Atiyah-Patodi-Singer $`\eta `$-invariant. The relation between $`\mu `$ and $`\nu `$ is similar to that between the $`\eta `$ and the Chern-Simons invariants: more precisely, when $`\mu `$ is defined, then for a CR structure $`J`$ one has
$$\nu (J)=3\mu (J)+\text{constant},$$
with the constant depending only of the underlying contact structure \[11, Theorem 1.3\]. Burns-Epsteinโs version of Miyaoka-Yau then reads, if $`M`$ is the boundary at infinity of a Kรคhler-Einstein $`N`$:
(91)
$$\chi (N)\frac{1}{3}\overline{c}_1(N)^2\mu (M),$$
with equality if the metric is complex hyperbolic; here $`\overline{c}_1`$ is a lift in $`H^2(N,M)`$ of $`c_1(N)`$.
A first important difference here is that our obstruction in Corollary 1.9 (filling by an ACH Einstein metric) is purely topological, whereas (91) involves a complex structure and a Kรคhler-Einstein metric.
Another important fact to be noticed, at least in the case when the quotient has no orbifold singularities, is that the obstructions obtained by both methods are different: if $`M`$ is a $`๐^1`$-bundle over the Riemann surface $`\mathrm{\Sigma }`$, then the $`\mu `$-invariant, being defined by a local formula, is multiplicative on finite coverings . Hence the values are
(92)
$$\mu =\frac{\chi ^2}{4d}\text{whereas}\nu =\frac{\chi ^2}{4d}d3.$$
Equation (91) implies that $`3\mu `$ must be an integer, i.e. $`\frac{3\chi ^2}{4d}`$ must belong to $``$, a condition that is weaker than Corollary 1.8, by a factor $`3`$.
Acknowledgements
. The authors are grateful to Yoshinobu Kamishima for useful conversations on the possible applications of $`\nu `$, and to Elisha Falbel for comments. M. H. thanks Emmanuel Royer for his help in computations of examples at a very early stage of this paper. Finally, we thank Nigel Hitchin for inventing the word โdiabaticโ.
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# Wave functions in the ๐_{7/2} shell, for educational purposes and ideas
## I Introduction
In this work, we present energy levels and wave functions for nuclei in the $`f_{7/2}`$ shell. There was a previous work by McCullen, Bayman and Zamick bmz63 ; mbz64 which was followed by a technical report that also contained energy levels and wave functions. In this compilation, as in the previous one, we take the two-body matrix elements from experiment. However, in 1964 some of the $`T=0`$ matrix elements were not known and guesses had to be made, some of which were incorrect; in this work, we use the correct spectrum. In the previous work, only the scandium and titanium nuclei, plus their cross conjugates and mirrors, were included; in the present compilation, we consider all the nuclei. For brevity, we do not include tables of wave functions for the cross conjugate or mirror nuclei, but we provide formulas that show how these can be obtained from their partners.
At present, one can perform complete $`fp`$ shell calculations, so what is the justification for calculations in such a truncated space? There is no question but that one needs a complete $`fp`$ space to get quantitative agreement with experiment, and even that is not always sufficientโintruder states can play an important role.
First of all, some properties do come out fairly well even in this small space. More importantly, we feel that just looking at the wave functions and performing simple calculations with them is a valuable educational experience. We provide examples of points of interest. Beyond that, some interesting symmetries arise that suggest new ideas. The large shell model programs do not provide wave functions. We feel that looking at these wave functions can lead to insights and surprises and indeed there will be some that the present authors have missed.
## II Coefficients of fractional parentage
In Tables 1 and 2, we list the one-particle coefficients of fractional parentage (cfpโs) used in this compilation. For a system of $`n`$ identical particles in a single $`j`$ shell with total angular momentum $`I`$, we have
$$\mathrm{\Psi }(j^nv\alpha IM)=\underset{v_1\alpha _1}{}[j^{n1}v_1\alpha _1J_1;j|\}j^nv\alpha I][\mathrm{\Psi }(j^{n1}v_1\alpha _1J_1)\mathrm{\Psi }(j)]^J_M.$$
(1)
In the above, $`\mathrm{\Psi }`$ is a totally antisymmetric wave function. The cfpโs are computed by the same method as was used by Bayman and Lande bl66 . We refer the reader to this work for details. It should, however, be pointed out that a different diagonalization routine is now being used to diagonalize the relevant matrices. Hence, some of the sets of cfpโs here can have an overall sign difference from the original BaymanโLande ones.
## III The interaction
We take the matrix elements from experiment, i.e., from the spectrum of <sup>42</sup>Sc. We make the association
$$(f_{7/2}^2)^J|V|(f_{7/2}^2)^J=E^{}(J),$$
(2)
where $`E^{}(J)`$ is the experimental excitation energy in <sup>42</sup>Sc of the lowest state of angular momentum $`J`$.
The values of $`E^{}(J)`$ in MeV are
| $`T=1`$ | | | $`T=0`$ | |
| --- | --- | --- | --- | --- |
| $`J=0`$ | 0.0000 | | $`J=1`$ | 0.6111 |
| $`J=2`$ | 1.5863 | | $`J=3`$ | 1.4904 |
| $`J=4`$ | 2.8153 | | $`J=5`$ | 1.5101 |
| $`J=6`$ | 3.2420 | | $`J=7`$ | 0.6163 |
Note that, for the $`(j^2)`$ configuration, the states with even $`J`$ have isospin $`T=1`$ and those with odd $`J`$ have isospin $`T=0`$.
It should be noted that the above energies differ from those of M.B.Z. mbz64 because in 1963 the $`T=0`$ states were not well known. The M.B.Z. values for $`T=1`$ were ($`0,1.509,2.998`$, and $`3.400`$), while for $`T=0`$ they were ($`1.035,2.248,1.958`$, and $`0.617`$).
## IV Calculation of the wave functions
For previous works on the $`f_{7/2}`$ shell, we advise the reader to consult Refs. bmz63 ; mbz64 ; gf63 ; g65 ; g66 ; kbo78 .
We present results for up to four protons and for any number of neutrons. With a charge independent interaction, such as the one we use here, this covers all cases because the nuclei not included are either mirror nuclei or cross conjugate nuclei. A mirror nucleus is one in which we change protons into neutrons and neutrons into protons. Assuming charge symmetry, the spectra of mirrors are identical. This result holds for multishell wave functions. For an explanation of cross conjugate nuclei, see Subsect. IV.1.
If the neutron number is four or less, we use the coefficients of fractional parentage of the BaymanโLande method previously described. For neutron number greater than four, we switch to neutron holes. Operationally, we perform the same calculation as before, but for the protonโneutron interaction we substitute the protonโneutron-hole interaction, which can be obtained by a Pandya relation
$$E(jj^1J)=\text{constant}\underset{J^{}}{}(2J^{}+1)\left\{\begin{array}{ccc}j& j& J^{}\\ j& j& J\end{array}\right\}E(jjJ^{}),$$
(3)
or, what is equivalent, from the calculated spectrum of <sup>48</sup>Sc. We do not change the protonโproton interaction and we note that the spectrum of two neutron holes is the same as that of two neutrons. The values of the protonโneutron-hole matrix elements, with the lowest state energy ($`J=6^+`$) set to zero for $`J=0,1,\mathrm{},7`$ are, respectively, (in MeV): 7.03763, 2.39516, 0.32755, 0.39826, 0.12027, 0.08821, 0.00000, 1.18812.
### IV.1 Cross conjugate nuclei
In the single-$`j`$-shell model, consider a nucleus with $`๐ซ`$ valence protons and $`๐ฉ`$ valence neutrons. We define the cross conjugate nucleus as one with $`๐ฉ`$ proton holes and $`๐ซ`$ neutron holes, i.e., ($`2j+1๐ฉ`$) protons and ($`2j+1๐ซ`$) neutrons. In the single-$`j`$-shell model, the spectra of these two nuclei are identical. If the wave function of the original nucleus is
$$\mathrm{\Psi }=D^I(J_P,J_N)\left[(j_\pi ^๐ซ)^{J_P}(j_\nu ^๐ฉ)^{J_N}\right]^I$$
(4)
and the wave function of the cross conjugate nucleus is
$$\mathrm{\Phi }=C^I(J_N,J_P)\left[(j_\pi ^{2j+1๐ฉ})^{J_N}(j_\nu ^{2j+1๐ซ})^{J_P}\right]^I,$$
(5)
then, for $`N4,Z4`$,
$$C^I(J_N,J_P)=D^I(J_P,J_N)(1)^{J_P+J_NI}.$$
(6)
For $`N=4,Z4`$ or $`N4,Z=4`$,
$$C^I(J_N,J_P)=D^I(J_P,J_N)(1)^{J_P+J_NI}(1)^{v_4/2},$$
(7)
where $`v_4`$ is the seniority of the four-particle system.
For $`N=Z=4`$, if we consider 4 protons and 4 neutron holes, we get
$$\mathrm{\Psi }=\underset{J_P,J_N}{}D^I(J_P,J_N)(1)^{v_P/2}\left[(j^4)^{J_Pv_P}(j^4)^{J_Nv_N}\right]^I.$$
(8)
The wave function of a mirror nucleus is
$`\mathrm{\Psi }`$ $`=`$ $`{\displaystyle D^I(J_P,J_N)\left[(j_\nu ^๐ซ)^{J_P}(j_\pi ^๐ฉ)^{J_N}\right]^I}`$ (9)
$`=`$ $`{\displaystyle D^I(J_P,J_N)(1)^{J_P+J_NI}\left[(j_\pi ^๐ฉ)^{J_N}(j_\nu ^๐ซ)^{J_P}\right]^I}.`$
Here is a list of nuclei and their cross conjugates:
| Nucleus | Cross conjugate | | Nucleus | Cross conjugate |
| --- | --- | --- | --- | --- |
| <sup>43</sup>Sc | <sup>53</sup>Fe | | <sup>44</sup>V | <sup>52</sup>Co |
| <sup>44</sup>Sc | <sup>52</sup>Mn | | <sup>45</sup>V | <sup>51</sup>Fe |
| <sup>45</sup>Sc | <sup>51</sup>Cr | | <sup>46</sup>V | <sup>50</sup>Mn |
| <sup>46</sup>Sc | <sup>50</sup>V | | <sup>47</sup>V | <sup>49</sup>Cr |
| <sup>47</sup>Sc | <sup>49</sup>Ti | | <sup>48</sup>V | <sup>48</sup>V |
| <sup>48</sup>Sc | <sup>48</sup>Sc | | | |
| <sup>43</sup>Ti | <sup>53</sup>Co | | <sup>45</sup>Cr | <sup>51</sup>Co |
| <sup>44</sup>Ti | <sup>52</sup>Fe | | <sup>46</sup>Cr | <sup>50</sup>Fe |
| <sup>45</sup>Ti | <sup>51</sup>Mn | | <sup>47</sup>Cr | <sup>49</sup>Mn |
| <sup>46</sup>Ti | <sup>50</sup>Cr | | <sup>48</sup>Cr | <sup>48</sup>Cr |
| <sup>47</sup>Ti | <sup>49</sup>V | | | |
| <sup>48</sup>Ti | <sup>48</sup>Ti | | | |
## V Format of tables
The wave functions for a system of $`p`$ protons and $`n`$ neutrons with total angular momentum $`I`$ are represented as column vectors
$$\mathrm{\Psi }^{I\alpha }=\underset{J_PJ_N}{}D^{I\alpha }(J_P,J_N)\left[(j^p)^{J_P}(j^n)^{J_N}\right]^I.$$
(10)
We list the excitation energies at the top; below we list $`J_P,J_N`$ and the corresponding $`D^{I\alpha }`$.
Note that, as we go down the column, we find the orthonormality condition
$$J_P,J_ND^{I\alpha }(J_P,J_N)D^{I\alpha ^{}}(J_P,J_N)=\delta _{\alpha \alpha ^{}}.$$
(11)
We also have along the row the completeness condition
$$\underset{\alpha }{}D^{I\alpha }(J_P,J_N)D^{I\alpha }(J_P^{},J_N^{})=\delta _{J_PJ_P^{}}\delta _{J_NJ_N^{}}.$$
(12)
Unless specified otherwise, the isospin of the states is the lowest possible one: $`T=|NZ|/2`$. The higher isospin states are specified. One can determine the isospin of the states by adding to the original interaction a term $`a+bt(1)t(2)`$. This will not change the wave functions or the relative excitation energies for states of a given isospin, but it will cause a $`T(T+1)`$ splitting of states of different isospins.
As an example, we consider the two $`I=7/2`$ states in <sup>43</sup>Sc with excitation energies 0.00000 and 4.14201 MeV. The wave functions are, respectively
| $`E=0.00000`$ | $`0.78776[j(j^2)^0]^{7/2}+0.56165[j(j^2)^2]^{7/2}`$ |
| --- | --- |
| | $`+0.22082[j(j^2)^4]^{7/2}+0.12340[j(j^2)^6]^{7/2}`$ |
| $`E=4.14201`$ | $`0.50000[j(j^2)^0]^{7/2}+0.37268[j(j^2)^2]^{7/2}`$ |
| --- | --- |
| | $`+0.50000[j(j^2)^4]^{7/2}+0.60093[j(j^2)^6]^{7/2}`$ |
The 4.142 MeV state is assigned isospin $`T=3/2`$. Note in this case that the $`D(j,J_N)`$โs are identical with the $`I=7/2`$ coefficients of fractional parentage for 3 identical particles in Table 1.
As another example, consider the $`I=1/2`$ lowest energy state ($`E^{}=2.68294`$ MeV) of <sup>47</sup>Cr. This nucleus is not in the tables, but we can get its wave function from the mirror nucleus <sup>47</sup>V using the prescription of Eq. (9). Thus, the wave function is
$`\mathrm{\Psi }`$ $`=`$ $`0.24622[2,3/2]0.04793[2,3/2]+0.62848[2,5/2]+0.25132[2,5/2]`$
$`+0.12809[4,7/2]+0.65385[4,7/2]0.01776[4,9/2]+0.09999[4,9/2]`$
$`0.00413[5,9/2]0.01591[5,11/2]0.15530[6,11/2]0.01738[8,11/2],`$
where the asterisk indicates seniority 4. Note that there is more ($`J_P=2`$, $`v=2`$) admixture in the wave function than there is ($`J_P=2`$, $`v=4`$). On the other hand, there is more ($`J_P=4`$, $`v=4`$) than ($`J_P=4`$, $`v=2`$) admixture.
## VI Justification
We feel that a good justification for publishing the single-$`j`$-shell wave functions is the educational value they contain.
There are some striking observations that the reader can ponder:
1. The even $`J`$ states in <sup>42</sup>Sc ($`J=0,2,4`$, and 6) have isospin 1, and the odd $`J`$ states ($`J=1,3,5`$, and 7) have isospin 0.
2. In the evenโeven Ti isotopes, there is one $`I=0`$ state with isospin $`T=T_{\text{min}}+2`$, where $`T_{\text{min}}=|NZ|/2`$. There are no $`I=0`$ states with $`T=T_{\text{min}}+1`$.
3. The $`D(J_P,J_N)`$โs for $`I=j`$ in <sup>43</sup>Sc are the same as those of $`I=0`$ in <sup>44</sup>Ti; and the excitation energies in <sup>43</sup>Sc ($`I=j`$) are half of those in <sup>44</sup>Ti ($`I=0`$).
4. Except for some overall phases, the $`D(J_P,J_N)`$โs for the Scandium states with isospin $`|NZ|/2+1`$ are the coefficients of fractional parentage $`[j^nJ_Nj|\}j^{n+1}Iv]`$.
5. In <sup>44</sup>Ti the isospin $`T=0`$ and $`T=2`$ states are such that $`D(J_P,J_N)=D(J_N,J_P)(1)^{J_P+J_NI}`$, while for $`T=1`$, $`D(J_P,J_N)=D(J_N,J_P)(1)^{J_P+J_N+1I}`$.
6. In self-conjugate <sup>48</sup>Ti, a given state has either $`D(J_P,J_N)=D(J_N,J_P)`$ or $`D(J_P,J_N)=D(J_N,J_P)`$. Thus, we can assign a signature quantum number to the states. Note that the lowest two $`6^+`$ states have opposite signatures; both have isospin $`T=2`$ and are nearly degenerate.
7. In <sup>46</sup>Ti $`I=0`$, the column vector for the unique $`T=3`$ state has $`D(J_P,J_N)`$ identical to the $`D(j,J_N)`$ for the $`T=5/2`$ state of <sup>45</sup>Sc with $`I=j`$. This implies an equality between two-particle cfpโs and one-particle cfpโs.
8. Note that, in <sup>48</sup>Cr, $`(1)^{(v_P+v_N)/2}`$ is a good quantum number, where $`v_P`$ and $`v_N`$ are the seniorities of the protons and neutrons. This leads to striking visual patterns in the wave functions.
9. Because of the small difference in excitation energy of the $`19/2^{}`$ and $`15/2^{}`$ states in <sup>43</sup>Sc (<sup>43</sup>Ti), the $`19/2^{}`$ state is isomeric. Likewise the $`12^+`$ and $`10^+`$ states in <sup>44</sup>Ti. However, in the cross conjugate nucleus <sup>52</sup>Fe, the $`12^+`$ comes below the $`10^+`$ and so it has a much bigger lifetime.
10. Note that the spectrum of <sup>48</sup>Sc (particleโhole) is almost the inverted spectrum of <sup>42</sup>Sc. With a $`QQ`$ interaction, the spectrum of <sup>48</sup>Sc would be precisely the upside down spectrum of <sup>42</sup>Sc.
11. Note that the single-$`j`$-shell model predicts the near fivefold degeneracy for the ground state of <sup>46</sup>Sc. The spins involved are $`I=2,3,4,5`$, and 6. In <sup>45</sup>Ti the calculation leads to a near degeneracy of the states with $`I=5/2^{}`$ and $`7/2^{}`$.
12. There is an interesting experimental systematic for the Sc isotopes. The ground state spins for <sup>42</sup>Sc, <sup>44</sup>Sc, <sup>46</sup>Sc, and <sup>48</sup>Sc are, respectively, $`I=0,2,4`$, and 6.
13. There is a dominance of $`J=0`$ and $`J=2`$ couplings in the lowest-lying states of the evenโeven nuclei. This lends some support to theoretical models which truncate to these two couplings, such as the interacting boson model.
14. In the low-lying states of systems with 4 neutrons or 4 protons, there is more admixture of $`I=4`$, $`v=4`$ than there is of $`I=4`$, $`v=2`$. This is in part due to the fact that in <sup>44</sup>Ca the $`I=4`$, $`v=4`$ state is slightly lower in energy than $`I=4`$, $`v=2`$. If $`I=4`$ is important, then seniority truncation will not work.
15. We note that the idea of taking matrix elements from experiment is present in the book of de Shalit and Talmi st63 . For the upper half of the $`f_{7/2}`$ shell, one might use the spectrum of <sup>54</sup>Co (holeโhole) rather than that of <sup>42</sup>Sc.
16. Note that, for the evenโeven and oddโeven nuclei considered here, the states of higher isospin are at high excitation energies. These states have a simpler structure than their lower isospin neighbours. Assuming the neutron has isospin $`+1/2`$ and the proton $`1/2`$, these higher isospin states in a nucleus $`(N,Z)`$ can be obtained by applying the isospin lowering operator to a state in the nucleus $`(N+1,Z1)`$.
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# Energetic disorder at the metal/organic semiconductor interface
## Abstract
The physics of organic semiconductors is dominated by the effects of energetic disorder. We show that image forces reduce the electrostatic component of the total energetic disorder near an interface with a metal electrode. Typically, the variance of energetic disorder is dramatically reduced at the first few layers of organic semiconductor molecules adjacent to the metal electrode. Implications for charge injection into organic semiconductors are discussed.
The past two decades have been characterized by dramatic advances in the performance of organic semiconductor devices, giving rise to the field known as organic electronics A . Light-emitting diodes B , thin film transistors C , and photovoltaic cells D , are examples of devices being developed based on organic semiconductors. Critical to the operation of all these devices is the process of charge injection from a metal electrode into the organic semiconductor. The efficiency of organic light emitting diodes, for example, is directly related to the ability of the contacts to supply the organic bulk with charge E . Despite the great technological importance of charge injection, the physics of this process remains poorly understood. This may be ascribed to the fact that transport in organic semiconductors is very different from that in their inorganic counterparts. In the former materials, transport takes place by hopping between localized electronic states, which are distributed in energy due to spatially correlated energetic disorder with the standard deviation $`\sigma 0.1`$ eV dip2 ; dpk .
Recently, it was recognized that energetic disorder in organic materials used in todayโs devices affects the injection efficiency conwell ; 1 ; 2 . First, disorder was shown to increase injection and, second, it was proposed as a major reason for the unusually weak temperature dependence of the injected current 3 ; 4 .
The injection properties of metal/organic interfaces depend on the properties of a thin organic layer directly contacting with the metal. It is well known that the structure of this interface layer is typically different from the bulk structure of the organic material. For this reason we may suspect that the variance $`\sigma _i^2`$ of the energetic disorder at the interface differs from the variance $`\sigma _b^2`$ of the disorder in the bulk of the organic material. In literature conwell ; 1 ; 2 , a calculation of the effect of energetic disorder on the injection has been carried out using bulk disorder parameters (basically, its variance $`\sigma _b^2`$). To some extent this could be explained by lack of any detailed knowledge of the structure of the interface layer. In addition, experimental data of temperature dependence of the injected current seem to supports the idea that $`\sigma _i\sigma _b`$ 3 . At the same time, it is well known that frequently a surface dipole layer is formed directly at the interface, providing an abrupt leap in carrier energy in the range of $`0.31`$ eV ishii . It is reasonable to assume that such a layer has some degree of disorder and, thus, induces additional energetic disorder in neighboring layers of organic materials 4 . The magnitude of this energetic disorder should decay while going away from the interface, so $`\sigma _i>\sigma _b`$, but this magnitude is completely unknown; in the calculations carried out in Ref. 4 very speculative parameters have been used to estimate $`\sigma _i`$.
We are going to show that this problem of the relative magnitude of $`\sigma _i`$ and $`\sigma _b`$ has an additional and quite remarkable twist, because in organic devices sandwiched between conducting electrodes the bulk disorder itself depends on the proximity to the electrode. A well known fact is that a significant part of the total energetic disorder in organic materials has an electrostatic origin. For polar materials this is dipolar disorder dip2 ; dpk ; dip1 , while for non-polar materials it is quadrupolar disorder q2 . Our major goal is to demonstrate that the electrostatic disorder at the vicinity of metal/organic interface differs from the bulk disorder far away from the electrode.
Indeed, the electrostatic energetic disorder in organic materials is directly proportional to the disorder in the spatial distribution of electrostatic potential, generated by randomly situated and oriented dipoles or quadrupoles. In organic layers sandwiched between conducting electrodes this spatial distribution must obey a boundary condition at the electrode surface: at this surface the potential should be a constant. Thus, at the electrode surface there is no energetic disorder at all, irrespectively to how disordered is the material in the organic layer. This means that the magnitude of the dipolar or quadrupolar disorder increases while going away from the interface, asymptotically reaching its bulk value. Here we assume that there is no significant increase of the dipolar or quadrupolar disorder directly at the interface (i.e., a disordered surface dipole layer is absent). Now we are going to support this general idea with the calculation of the variance of the dipolar disorder in the vicinity of a conducting electrode.
Let us consider the simplest model of a rigid disordered polar organic material where the randomly oriented (and orientationally uncorrelated) point dipoles occupy the sites of a simple cubic lattice with the lattice scale $`a`$ dip1 ; dip2 . We consider the vicinity of a metal electrode located at $`z=0`$, so the lattice occupies the half-space $`z>0`$ with the first lattice layer having distance $`a/2`$ from the electrode plane. Charge carrier energy at any site $`m`$ is the sum
$$U(\stackrel{}{r}_m)=e\underset{nm}{}\varphi (\stackrel{}{r}_m,\stackrel{}{r}_n)$$
(1)
where $`\varphi (\stackrel{}{r},\stackrel{}{r}_n)`$ is the electrostatic potential, generated by the dipole, located at the site $`n`$. The variance of the disorder is
$$\sigma ^2(\stackrel{}{r}_m)=U_m^2=e^2\underset{n,lm}{}\xi _n\xi _l\varphi (\stackrel{}{r}_m,\stackrel{}{r}_n)\varphi (\stackrel{}{r}_m,\stackrel{}{r}_l),$$
(2)
where the angular brackets denote the average over positions and orientations of dipoles, and the variable $`\xi _n`$ equals to 1 if a dipole is located at the site $`n`$ and 0 otherwise (note that $`U_m=0`$). A spatial average gives
$$\xi _n\xi _l=c\delta _{nl}+c^2\left(1\delta _{nl}\right),$$
(3)
where $`c`$ is the fraction of sites occupied by dipoles, and taking into account that the average over dipole orientations gives $`\varphi (\stackrel{}{r}_m,\stackrel{}{r}_n)=0`$, we obtain
$$\sigma ^2(\stackrel{}{r}_m)=e^2c\underset{nm}{}\varphi ^2(\stackrel{}{r}_m,\stackrel{}{r}_n).$$
(4)
From this point the angular brackets denote the orientational average only. In the case of an infinite medium without any electrodes
$$\varphi (\stackrel{}{r},\stackrel{}{r}_n)=\frac{\stackrel{}{p}_n\left(\stackrel{}{r}\stackrel{}{r}_n\right)}{\epsilon \left|\stackrel{}{r}\stackrel{}{r}_n\right|^3},$$
(5)
where $`\epsilon `$ is the dielectric constant of the medium and $`\stackrel{}{p}_n`$ is the dipole moment. In the case of semi-infinite medium bounded by an electrode, a boundary condition $`\varphi =0`$ at $`z=0`$ has to be imposed (we choose the arbitrary constant to be zero). As a result, the source function $`\varphi (\stackrel{}{r},\stackrel{}{r}_n)`$ includes a contribution from the image dipole $`\stackrel{}{p}_n^\text{i}=(p_{nx},p_{ny},p_{nz})`$ located at $`\stackrel{}{r}_n^\text{i}=(x_n,y_n,z_n)`$. an average over dipole orientations gives $`p_{ni}p_{nj}=\frac{1}{3}p^2\delta _{ij}`$ (with the obvious modification for $`p_{ni}p_{nj}^\text{i}`$), and finally
$`\sigma ^2(z)={\displaystyle \frac{e^2p^2c}{3\epsilon ^2}}{\displaystyle \underset{z_n>0}{}}[{\displaystyle \frac{1}{\left|\stackrel{}{r}_n\stackrel{}{z}\right|^4}}+{\displaystyle \frac{1}{\left|\stackrel{}{r}_n+\stackrel{}{z}\right|^4}}`$ (6)
$`2{\displaystyle \frac{r_n^2z^2}{\left|\stackrel{}{r}_n\stackrel{}{z}\right|^3\left|\stackrel{}{r}_n+\stackrel{}{z}\right|^3}}],`$
here the vector $`\stackrel{}{z}=(0,0,z)`$ and $`\sigma `$ does not depend on $`x`$ and $`y`$. The lattice site with $`\stackrel{}{r}_n=\stackrel{}{z}`$ is excluded from the sum (6). Note that $`\sigma (0)=0`$, as it should be, because if the electrostatic potential is a constant for $`z=0`$, then there is no electrostatic disorder at this plane, no matter how many dipoles occupy the half-space $`z>0`$. Far away from the electrode dip2
$$\sigma ^2(\mathrm{})=\sigma _b^2=\frac{e^2p^2c}{3\epsilon ^2}\underset{r_n0}{}\frac{1}{r_n^4}5.51\frac{e^2p^2c}{\epsilon ^2a^4}.$$
(7)
We can perform an approximate analytical summation in Eq. (6) if we provide the alternative expression for the source function $`\varphi (\stackrel{}{r},\stackrel{}{r}_0)`$, in close analogy to the method, used in Ref. synthmet . The source function for the point dipole, located at $`\stackrel{}{r}_0`$, is the solution of the Poisson equation
$$\mathrm{\Delta }\varphi =\frac{4\pi }{\epsilon }\stackrel{}{p}_{\stackrel{}{r}_0}\delta (\stackrel{}{r}\stackrel{}{r}_0)$$
(8)
and, hence, is proportional to the gradient of the Green function $`G(\stackrel{}{r},\stackrel{}{r}_0)`$ of the Laplace operator with the zero boundary condition at $`z=0`$
$$\varphi (\stackrel{}{r},\stackrel{}{r}_0)=\frac{4\pi }{\epsilon }\stackrel{}{p}_{\stackrel{}{r}_0}G(\stackrel{}{r},\stackrel{}{r}_0).$$
(9)
Replacing summation with integration in Eq. (4) we have
$$\sigma ^2(z)\frac{16\pi ^2e^2p^2c}{3\epsilon ^2a^3}_{z^{}>0}๐\stackrel{}{r^{}}\left[_\stackrel{}{r^{}}G(\stackrel{}{z},\stackrel{}{r^{}})\right]^2$$
(10)
where the Green function has the form synthmet
$`G(\stackrel{}{r},\stackrel{}{r^{}})={\displaystyle \frac{1}{4\pi ^2}}{\displaystyle ๐\stackrel{}{k}e^{i\stackrel{}{k}(\stackrel{}{\rho }\stackrel{}{\rho }^{})}G_k(z,z^{})},`$ (11)
$`G_k(z,z^{})={\displaystyle \frac{\mathrm{sinh}kz_{}}{k}}\mathrm{exp}(kz_+),`$
$`z_+=\mathrm{max}(z,z^{}),z_{}=\mathrm{min}(z,z^{}),`$
and $`\stackrel{}{\rho }=(x,y)`$ and $`\stackrel{}{k}`$ are two-dimensional vectors. Performing integration over $`\stackrel{}{\rho }^{}`$ in Eq. (10) we obtain
$`\sigma ^2(z){\displaystyle \frac{8\pi e^2p^2c}{3\epsilon ^2a^3}}{\displaystyle _0^{\mathrm{}}}๐kk{\displaystyle _0^{\mathrm{}}}๐z^{}\left[k^2G_k^2(z,z^{})+\left({\displaystyle \frac{dG_k}{dz^{}}}\right)^2\right]`$ (12)
$`{\displaystyle \frac{4\pi e^2p^2c}{3\epsilon ^2a^3}}{\displaystyle _0^{1/a_0}}๐k\left(1e^{2kz}\right)=\sigma _b^2\left[1{\displaystyle \frac{a_0}{2z}}\left(1e^{2z/a_0}\right)\right],\sigma _b^2={\displaystyle \frac{4\pi e^2p^2c}{3\epsilon ^2a^3a_0}}.`$
Here a cut-off at $`k1/a_0`$ with $`a_0a`$ has been introduced to remove the divergence at $`k\mathrm{}`$. This cut-off is equivalent to the exclusion of the site with $`\stackrel{}{r}_n=\stackrel{}{z}`$ in Eq. (6). We did not introduce a similar short range cut-off in the integral over $`z^{}`$ in Eq. (12) because this integral converges and the possible correction does not change the result in a qualitative way. To obtain an agreement between the bulk $`\sigma _b`$ in Eq. (12) and the corresponding exact value for the lattice model in Eq. (7) we have to set $`a_00.76a`$ SPIE97 . This choice of $`a_0`$ leads to a remarkably good agreement between the approximate analytic expression (12) and the result of the direct summation according to Eq. (6) in the whole range of distance from the interface (see Fig. 1).
Note that for the transport sites situated within a distance of $`56`$ lattice sites to the interface, the amplitude of energetic disorder is significantly decreased in comparison to its bulk value. Yet this very thin layer is of crucial importance for injection in organic devices. We anticipate that the reduction of $`\sigma `$ should lead to a stronger temperature dependence for the injected current density in comparison with the treatment with $`\sigma (z)=\text{const}`$, provided in the paper by Arkhipov et al. 1 . According to this model, the injected current density is
$$J_d^{\mathrm{}}๐z_0\mathrm{exp}(2\gamma z_0)w_{\text{esc}}(z_0)_{\mathrm{}}^{\mathrm{}}๐U^{}\text{Bol}(U^{})g\left[U_0(z_0)U^{}\right],U_0(z_0)=\mathrm{\Delta }\frac{e^2}{4\epsilon z_0}eEz_0$$
(13)
where $`\text{Bol}(U)=\mathrm{exp}(U/kT)`$ for $`U>0`$ and $`\text{Bol}(U)=1`$ otherwise, and $`w_{\text{esc}}(z)`$ is the probability for a carrier to avoid the surface recombination
$$w_{\text{esc}}(z_0)=\frac{_d^{z_0}๐z\mathrm{exp}\left[\frac{U_0(z)}{kT}\right]}{_d^{\mathrm{}}๐z\mathrm{exp}\left[\frac{U_0(z)}{kT}\right]}.$$
(14)
Here $`E`$ is the applied electric field, $`g(U)`$ is the density of states in the organic material (a Gaussian density of states is usually assumed), $`\gamma `$ is the inverse localization radius, $`\mathrm{\Delta }`$ is the interface barrier, and $`d`$ is the minimal distance, separating the electrode surface and the first layer of the organic material. A natural generalization of the Eq. (13) to our case is straightforward: we have to let the density of states $`g`$ depend on $`z`$ through the Eq. (12). We performed the calculation using parameters provided in Ref. 3 for the injection of holes from the Ag electrode into poly-dialkoxy-p-phenylene vinylene: $`\mathrm{\Delta }=0.95`$ eV, $`\gamma =0.33`$ ร
<sup>-1</sup>, $`E=5\times 10^5`$ V/cm, $`\sigma _b=0.11`$ eV, and $`d=12`$ ร
(it was assumed that $`d=a`$). Note that all these parameters were used in Ref. 3 for the comparison between the experimental data and Eq. (13) not as fitting parameters, but have been measured independently. The result of the calculation is shown in Fig. 2. A transformation of the curve occurs as anticipated: because of the smaller $`\sigma `$ at the interface, the temperature dependence of the injected current becomes stronger and does not agree with experimental data anymore. In fact, the agreement between the experimental points and the curve, calculated using Eq. (13) for $`\sigma (z)=\text{const}`$, is not as perfect as it appears to be, because we have to expect that $`d<a`$ is a better choice for the minimal distance to the electrode ($`a`$ is the intersite separation). For $`d<a`$ the curve, calculated by Eq. (13) for $`\sigma (z)=\text{const}`$, goes up (see Fig. 2, the upper solid curve), and the agreement worsens. Additionally, small distances to the electrode are very important in the integral (13), so the use of the macroscopic $`\epsilon `$ in Eq. (13) is dubious. Again, the decrease in $`\epsilon `$ moves the curve for $`\sigma (z)=\text{const}`$ in Fig. 2 upwards.
All these reservations nothwithstanding, if we believe that the model of Arkhipov et al. 1 is valid, then the significant discrepancy between the lower broken line and the experimental points in Fig. 2 seems to be an indication of the need to introduce an additional disorder (with $`\sigma 0.1`$ eV) at the interface. The possible source of this additional disorder could be the surface dipole layer. From this point of view, the experimental results, provided in Ref. 3 and connected to our consideration, in fact support the idea of a disordered surface dipole layer: we clearly need additional disorder at the interface to compensate the decrease in the electrostatic disorder, provided by the bulk molecules.
Possible generalizations (taking into account the possible additional spatial disorder at the interface, roughness of the metal/organic interface) do not change our main conclusion that the electrostatic part of the energetic disorder in organic materials, provided by molecules in the bulk of the material, is significantly suppressed at the interface.
If a disordered surface dipolar layer is indeed formed at the interface, the picture suggested in this paper has to be modified. In this case the magnitude of the total energetic disorder could decrease with the increase of the distance to the electrode, or it could still increase, depending on the relative amplitudes of the bulk and surface contributions.
In conclusion, we have shown that the energetic disorder in organic semiconductors decreases dramatically in the neighborhood of a metal electrode. The ramification of this study is that disorder parameters derived from bulk measurements do not describe first few layers near the metal. As a result, simple models that predict enhancement of charge injection in organic semiconductors due to presence of disorder need to be reexamined.
This work was supported by the ISTC grant 2207 and RFBR grants 02-03-33052 and 03-03-33067. The research described in this publication was made possible in part by Award No. RE2-2524-MO-03 of the U.S. Civilian Research & Development Foundation for the Independent States of the Former Soviet Union (CRDF).
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# A moduli approach to quadratic โ-curves realizing projective mod ๐ Galois representations
## 1 Introduction
Throughout we fix an odd prime $`p`$ and an algebraic closure $`\overline{}`$ of $``$. For a subfield $`L`$ of $`\overline{}`$, we denote by $`\mathrm{G}_L`$ the absolute Galois group $`\mathrm{Gal}(\overline{}/L)`$. We say that a non-CM elliptic curve defined over a quadratic field $`k`$ is a *$``$-curve of degree $`N`$* if there is a cyclic isogeny of degree $`N`$ from the curve to its Galois conjugate. The $`p`$-torsion of a $``$-curve $`E_{/k}`$ of degree $`N`$ prime to $`p`$ gives rise to a Galois representation
$$\varrho _E:\mathrm{G}_{}\mathrm{PGL}_2(๐ฝ_p)$$
whose conjugacy class is an invariant of the isomorphism class of $`E`$ and whose restriction to $`\mathrm{G}_k`$ is the projective representation obtained from the usual Galois action on the $`p`$-torsion points of the curve. The procedure to obtain $`\varrho _E`$ is detailed in Section 2. The determinant of $`\varrho _E`$ draws off two different cases, which we call *cyclotomic* and *non-cyclotomic*, and which correspond to $`N`$ being a square mod $`p`$ or not, respectively. These two cases rule most of the structure and contents of the rest of sections.
The situation just described raises the following inverse problem: find the $``$-curves of degree $`N`$ realizing a given projective mod $`p`$ Galois representation
$$\varrho :\mathrm{G}_{}\mathrm{PGL}_2(๐ฝ_p),$$
namely those whose $`p`$-torsion points give rise to $`\varrho `$ in the above sense. The aim of this paper is to explain in detail how to produce the moduli spaces whose rational points yield the solutions to this problem. Henceforth the term *rational* stands for *$``$-rational*.
The moduli spaces that we provide are either twists of the modular curve $`X(N,p)`$ obtained as the fiber product of the curves $`X_0(N)`$ and $`X(p)`$, in the non-cyclotomic case, or twists of a certain Atkin-Lehner quotient $`X^+(N,p)`$ in the cyclotomic case. In Section 3 we analyze the structure of the subgroup $`๐ฒ(N,p)`$ of automorphisms on $`X(N,p)`$ extending the group generated by the Atkin-Lehner involution $`w_N`$ on $`X_0(N)`$. In Section 4 we fix a suitable rational model for $`X(N,p)`$ and then describe the Galois action on $`๐ฒ(N,p)`$. In order to do this, we need first to study the action of $`๐ฒ(N,p)`$ on the non-cuspidal points of the curve. The last two sections explain how to obtain the twisted curves whose non-cuspidal non-CM rational points give the $``$-curves of degree $`N`$ realizing $`\varrho `$. Section 5 is devoted to the cyclotomic case and Section 6 to the non-cyclotomic case. They also include some finiteness results obtained from Faltingsโ theorem and from some genus computations performed in Section 3.
Our moduli approach turns out to be quite effective in the non-cyclotomic case, since the quadratic field of definition for the possible $``$-curves realizing the representation $`\varrho `$ is uniquely determined and, for every fixed degree $`N`$โ, we just need one twist $`X(N,p)_\varrho `$ . For an explicit application we refer to \[FGL\], where a plane quartic model is provided for the genus-three case $`X(5,3)_\varrho `$ .
In the cyclotomic case, one should instead consider two twists $`X^+(N,p)_\varrho `$ , $`X^+(N,p)_\varrho ^{}`$ whose rational points include the cyclic isogenies of degree $`N`$ between elliptic curves over $``$ realizing $`\varrho `$. One may also approach the problem by adding a given quadratic field $`k`$ as extra data: the $``$-curves of degree $`N`$ defined over $`k`$ realizing $`\varrho `$ are given by the non-cuspidal non-CM rational points on two other twisted curves $`X(N,p)_{\varrho ,k}`$ , $`X(N,p)_{\varrho ,k}^{}`$ .
## 2 Projective mod $`p`$ Galois representations realized by $`p`$-admissible $``$-curves
The aim of this section is to review the construction of the representation
$$\varrho _E:\mathrm{G}_{}\mathrm{PGL}_2(๐ฝ_p)$$
attached to a ($`p`$-admissible) $``$-curve $`E`$ and to compute its determinant. Up to some minor points, the section is mostly a reformulation of known facts that can mainly be found in \[ES01\] and go back to \[Rib92\]. The particular case of quadratic $``$-curves is written down in \[Ser92\] using the ideas of \[Shi78\].
Let $`E`$ be a *$``$-curve*. By this we mean a non-CM elliptic curve defined over a number field $`L`$ and with an isogeny
$$\lambda _\sigma :^\sigma EE$$
for every $`\sigma `$ in $`\mathrm{G}_{}`$. Without loss of generality, we always take $`\lambda _\sigma `$ equal to $`\lambda _\tau `$ whenever $`\sigma `$ and $`\tau `$ restrict to the same embedding of $`L`$ into $`\overline{}`$, and one might also assume the isogenies $`\lambda _\sigma `$ to be cyclic. We suppose here that the $``$-curve $`E`$ is *$`p`$-admissible*, namely that the isogenies $`\lambda _\sigma `$ can be chosen so that $`p`$ does not divide the degree of any of them.
For an isogeny $`\phi :E^{}E`$, let us write $`\phi ^1`$ for the element $`(\mathrm{deg}\phi )^1\widehat{\phi }`$ in $`\mathrm{Hom}(E,E^{})`$, where $`\widehat{\phi }`$ is the dual isogeny of $`\phi `$. Since $`E`$ has no CM, any isogeny $`E^{}E`$ differs from $`\phi `$ by a rational number. Thus, the $`2`$-cocycle of $`\mathrm{G}_{}`$
$$c_E:(\sigma ,\tau )\lambda _\sigma {}_{}{}^{\sigma }\lambda _{\tau }^{}\lambda _{\sigma \tau }^1$$
takes values in $`^{}`$โ. Let $`\alpha `$ be a *splitting map* for the $`2`$-cocycle $`c_E`$ viewed inside the trivial cohomology group $`H^2(\mathrm{G}_{},\overline{}^{})`$, that is, a continuous map $`\mathrm{G}_{}\overline{}^{}`$ satisfying
$$\lambda _\sigma {}_{}{}^{\sigma }\lambda _{\tau }^{}\lambda _{\sigma \tau }^1=\alpha (\sigma )\alpha (\tau )\alpha (\sigma \tau )^1$$
for all $`\sigma ,\tau `$ in $`\mathrm{G}_{}`$. By taking degrees, one deduces that the map $`\sigma \alpha (\sigma )^2/\mathrm{deg}\lambda _\sigma `$ is a Galois character. In particular, the values taken by $`\alpha `$ are algebraic integers prime to $`p`$. So there exist a finite extension $`๐ฝ_\alpha `$ of $`๐ฝ_p`$ and a mod $`p`$ reduction map $`\stackrel{~}{\alpha }:\mathrm{G}_{}๐ฝ_\alpha ^{}`$ obtained from a fixed embedding of $`\overline{}`$ into a fixed algebraic closure $`\overline{}_p`$ of $`_p`$.
Consider now the $`๐ฝ_\alpha `$โโlinear action of $`\mathrm{G}_{}`$ on $`๐ฝ_\alpha _{๐ฝ_p}E[p]`$ given by
$$(\sigma ,\mathrm{\hspace{0.17em}1}P)\stackrel{~}{\alpha }(\sigma )^1\lambda _\sigma ({}_{}{}^{\sigma }P).$$
By means of the choice of a basis for the $`๐ฝ_p`$โโmodule $`E[p]`$, this action produces a linear representation
$$\rho _{E,\alpha }:\mathrm{G}_{}๐ฝ_\alpha ^{}\mathrm{GL}_2(๐ฝ_p)$$
defined up to conjugation by matrices in $`\mathrm{GL}_2(๐ฝ_p)`$. The corresponding projective Galois representation $`\varrho _E`$ is actually given by the induced action
$$(\sigma ,C)\lambda _\sigma (^\sigma C)$$
on the projective line
$$\left(E[p]\right)=\left\{CE[p]\right|C๐ฝ_p\}.$$
This projective representation $`\varrho _E`$ depends on neither the $`p`$-admissible system of isogenies $`\lambda _\sigma `$ nor the splitting map $`\alpha `$. Further, the following proposition shows that $`\varrho _E`$ is an invariant of the *$`p`$-admissible isogeny class* of $`E`$.
###### Proposition 2.1
Let $`E^{}`$ be an elliptic curve over $`\overline{}`$ and $`\phi :E^{}E`$ be an isogeny of degree prime to $`p`$. Then $`\varrho _E^{}=\varrho _E`$.
*Proof.* Let $`\widehat{\phi }`$ be the dual isogeny of $`\phi `$ and let $`\lambda _\sigma `$ and $`\alpha `$ be as before. Consider the $`2`$-cocycle $`c_E^{}`$ attached to the $`p`$-admissible system of isogenies $`\widehat{\phi }\lambda _\sigma {}_{}{}^{\sigma }\phi `$ for the $``$-curve $`E^{}`$โ. Then $`\alpha \mathrm{deg}\phi `$ is a splitting map for $`c_E^{}`$ whose reduction mod $`p`$ takes values in the same finite field $`๐ฝ`$ as $`\stackrel{~}{\alpha }`$. The isomorphism $`E^{}[p]E[p]`$ induced by $`\phi `$ extends naturally to an isomorphism $`๐ฝE^{}[p]๐ฝE[p]`$ that is compatible with the corresponding $`๐ฝ`$โlinear actions of $`\mathrm{G}_{}`$. So $`\rho _{E,\alpha }`$ and $`\rho _{E^{},\alpha \mathrm{deg}\phi }`$ are conjugated by a matrix in $`\mathrm{GL}_2(๐ฝ_p)`$, and the result follows. $`\mathrm{}`$
###### Remark 2.2
The (conjugacy class of the) representation $`\rho _{E,\alpha }`$ is the linear mod $`p`$ representation obtained from the Galois action on the abelian variety of $`\mathrm{GL}_2`$-type attached in \[Rib92\] to the $``$-curve $`E`$ and the splitting map $`\alpha `$. Moreover, any lifting of $`\varrho _E`$ into $`\mathrm{GL}_2(\overline{๐ฝ}_p)`$ is of the form $`\rho _{E,\alpha }`$ for some splitting map $`\alpha `$ for $`c_E`$.
Note that the restriction of $`\varrho _E`$ to $`\mathrm{G}_L`$ is the projective representation
$$\overline{\rho }_E:\mathrm{G}_L\mathrm{PGL}_2(๐ฝ_p)$$
obtained from the usual Galois action on the $`p`$-torsion points of $`E`$. In terms of number fields, this provides the fixed field of $`\varrho _E`$ with the following property: its composite with $`L`$ is the splitting field of the modular polynomial $`\mathrm{\Phi }_p(j_E;X)`$ over $`L`$, where $`j_E`$ stands for the $`j`$-invariant of the elliptic curve $`E`$. Whenever $`L`$ is normal over $``$ and $`\overline{\rho }_E`$ is surjective, this property singles out the fixed field of $`\varrho _E`$ among all Galois extensions of $``$ with group $`\mathrm{PGL}_2(๐ฝ_p)`$.
We recall that the determinant of $`\overline{\rho }_E`$ is the restriction to $`\mathrm{G}_L`$ of the quadratic Galois character
$$\epsilon :\mathrm{G}_{}๐ฝ_p^{}/๐ฝ_{p}^{}{}_{}{}^{2}\{\pm 1\}$$
obtained from the mod $`p`$ cyclotomic character $`\chi `$. The fixed field of $`\epsilon `$ is the only quadratic field $`k_p=(\sqrt{\pm p})`$ inside the $`p`$-th cyclotomic extension of $``$. Let us now show that the projective representation $`\varrho _E`$ is odd by first computing the determinant of a lifting.
###### Proposition 2.3
The determinant of $`\rho _{E,\alpha }`$ is the product of the mod $`p`$ cyclotomic character $`\chi `$ and the character $`\mathrm{G}_{}๐ฝ_\alpha ^{}`$ defined by $`\sigma \mathrm{deg}\lambda _\sigma /\stackrel{~}{\alpha }(\sigma )^2`$โ.
*Proof.* By virtue of the properties of the Weil pairing $`,_{E,p}`$ , the equalities
$$\lambda _\sigma (^\sigma P),\lambda _\sigma (^\sigma Q)_{E,p}=^\sigma P,\widehat{\lambda }_\sigma \lambda _\sigma (^\sigma Q)_{{}_{}{}^{\sigma }E,p}=\left(P,Q_{E,p}\right)^{\chi (\sigma )\mathrm{deg}\lambda _\sigma }$$
hold for any two points $`P,Q`$ in $`E[p]`$ and any $`\sigma `$ in $`\mathrm{G}_{}`$. Whenever $`[P,Q]`$ is a basis of $`E[p]`$, so is $`[\lambda _\sigma (^\sigma P),\lambda _\sigma (^\sigma Q)]`$. Moreover, the left-hand term in the above equalities is the power of $`P,Q_{E,p}`$ to the determinant of the basis change. Therefore, the result follows from the definition of $`\rho _{E,\alpha }`$ . $`\mathrm{}`$
###### Corollary 2.4
The Galois representation $`\varrho _E`$ is odd.
*Proof.* Since $`\varrho _E`$ does not depend on the $`p`$-admissible system of isogenies $`\lambda _\sigma `$ chosen, we can take $`\lambda _\sigma `$ as the identity for all $`\sigma `$ in $`\mathrm{G}_L`$. Fix a splitting map $`\alpha `$ for the $`2`$-cocycle $`c_E`$ obtained from the isogenies $`\lambda _\sigma `$. Note that the restriction of $`\alpha `$ to $`\mathrm{G}_L`$ is then a Galois character. Write $`\varsigma `$ for the complex conjugation in $`\mathrm{G}_{}`$ obtained by fixing an embedding $`\overline{}`$. We must prove the equality $`det\rho _{E,\alpha }(\varsigma )=1`$. Consider the isogeny $`\lambda _\varsigma {}_{}{}^{\varsigma }\lambda _{\varsigma }^{}:EE`$. It is given, on the one hand, by multiplication by $`\alpha (\varsigma )^2`$ and, on the other hand, by multiplication by $`\pm \mathrm{deg}\lambda _\varsigma `$. All we have to do is to pin down the latter sign. As a complex elliptic curve, $`E`$ is isomorphic to $`E_z=/(+z)`$ for some $`z`$ in the complex upper-half plane $``$. Through this isomorphism, $`\lambda _\varsigma {}_{}{}^{\varsigma }\lambda _{\varsigma }^{}`$ translates into the isogeny $`E_zE_z`$ induced by multiplication by $`\delta {}_{}{}^{\varsigma }\delta `$ for some $`\delta `$ in $`^{}`$โ. So the above sign is positive and thus $`\stackrel{~}{\alpha }(\varsigma )^2=\mathrm{deg}\lambda _\varsigma `$ in $`๐ฝ_p`$. Since $`\chi (\varsigma )=1`$, the result follows from Proposition 2.3. $`\mathrm{}`$
Let $`deg:\mathrm{G}_{}^{}/_{}^{}{}_{}{}^{2}`$ be the degree character induced by any $`p`$-admissible system of isogenies $`\lambda _\sigma :{}_{}{}^{\sigma }EE`$. Then, consider the mod $`p`$ degree character
$$deg_p:\mathrm{G}_{}๐ฝ_p^{}/๐ฝ_{p}^{}{}_{}{}^{2}\{\pm 1\}$$
obtained from $`deg`$ by composition with the natural map $`^{}/_{}^{}{}_{}{}^{2}_p^{}/_{p}^{}{}_{}{}^{2}`$โ. The following statement is a straightforward consequence of Proposition 2.3.
###### Corollary 2.5
The determinant of $`\varrho _E`$ is the product $`\epsilon deg_p`$.
###### Remark 2.6
If the map $`deg`$ is not trivial, its fixed field $`K_{deg}`$ is a composite of quadratic fields $`(\sqrt{a_1}),\mathrm{},(\sqrt{a_m})`$, where $`2^m`$ is the degree of $`K_{deg}`$ over $``$. For every $`l=1,\mathrm{},m`$ , take $`\sigma _l`$ in $`\mathrm{G}_{}`$ restricting to the non-trivial automorphism of $`K_{deg}`$ that fixes $`\sqrt{a_h}`$ for $`hl`$. Then, the map $`deg_p`$ is the product of the quadratic Galois characters attached to the extensions $`(\sqrt{a_l})`$ for which $`\mathrm{deg}\lambda _{\sigma _l}`$ is not a square mod $`p`$.
We say that a projective mod $`p`$ Galois representation
$$\varrho :\mathrm{G}_{}\mathrm{PGL}_2(๐ฝ_p)$$
is *realized by* a ($`p`$-admissible) $``$-curve $`E`$ if $`\varrho _E=\varrho `$, where this equality makes only sense up to conjugation in $`\mathrm{PGL}_2(๐ฝ_p)`$. The rest of sections are devoted to the particular case of quadratic $``$-curves. Assume $`\varrho `$ to be realized by a $`p`$-admissible $``$-curve of degree $`N`$โ, that is, by a non-CM elliptic curve defined over a quadratic field and with a cyclic isogeny to its Galois conjugate of degree $`N`$ prime to $`p`$. From Corollary 2.5 and Remark 2.6, $`\varrho `$ has determinant $`\epsilon `$ if and only if $`N`$ is a square mod $`p`$, and otherwise any $``$-curves of degree $`N`$ realizing $`\varrho `$ must be defined over the fixed field of the quadratic character $`\epsilon det\varrho `$. We refer to the first case ($`N`$ square mod $`p`$) as the *cyclotomic case*, and to the second one ($`N`$ non-square mod $`p`$) as the *non-cyclotomic case*.
## 3 Automorphisms of the modular curve $`X(N,p)`$
Let $`N>1`$ be an integer prime to $`p`$. Let $`X_0(N)`$, $`X(p)`$ and $`X(1)`$ be the modular curves attached to the congruence subgroups $`\mathrm{\Gamma }_0(N)`$, $`\mathrm{\Gamma }(p)`$ and $`\mathrm{SL}_2()`$, respectively. We denote by $`X(N,p)`$ the modular curve attached to the congruence subgroup $`\mathrm{\Gamma }_0(N)\mathrm{\Gamma }(p)`$, namely the fiber product of $`X_0(N)`$ and $`X(p)`$ over $`X(1)`$ :
The aim of this section is to introduce a certain group $`๐ฒ(N,p)`$ of automorphisms on $`X(N,p)`$. We also compute the genus of this curve.
As a complex curve, $`X(N,p)`$ is a Galois covering of $`X_0(N)`$ with group $`๐ข(N,p)`$ given by the quotient $`\mathrm{\Gamma }_0(N)/\pm \mathrm{\Gamma }_0(N)\mathrm{\Gamma }(p)`$. Since the mod $`p`$ reduction map $`\mathrm{SL}_2()\mathrm{SL}_2(๐ฝ_p)`$ induces the exact sequence
$$1\pm \mathrm{\Gamma }_0(N)\mathrm{\Gamma }(p)\mathrm{\Gamma }_0(N)\mathrm{PSL}_2(๐ฝ_p)1,$$
there is a canonical isomorphism
$$๐ข(N,p)\mathrm{PSL}_2(๐ฝ_p).$$
We recall that $`๐ข(N,p)`$ consists of the automorphisms $`g`$ on $`X(N,p)`$ for which the following diagram commutes:
Let $`w_N`$ be the Atkin-Lehner involution on $`X_0(N)`$ and denote by $`X^+(N)`$ the corresponding quotient. For any integers $`a,b,c,d`$ satisfying $`adNbcp^2=1`$ and $`d\pm 1(\mathrm{mod}p)`$, the action of the matrix
$$\left(\begin{array}{cc}aN& bp\\ cpN& dN\end{array}\right)$$
on the complex upper-half plane $``$ defines an automorphism $`\vartheta `$ on $`X(N,p)`$ *extending* $`w_N`$, namely making the following diagram commutative:
Indeed, one can check that the above matrix lies in the normalizer of $`\mathrm{\Gamma }_0(N)\mathrm{\Gamma }(p)`$ inside $`\mathrm{PSL}_2()`$. Hence, the covering $`X(N,p)X^+(N)`$ has as many automorphisms as its degree, which means that it is a Galois covering. Let $`๐ฒ(N,p)`$ denote its automorphism group:
The group $`๐ฒ(N,p)`$ contains $`๐ข(N,p)`$ as a subgroup of index two whose complement consists of the automorphisms on $`X(N,p)`$ extending $`w_N`$.
###### Proposition 3.1
The group $`๐ข(N,p)`$ is a direct factor of $`๐ฒ(N,p)`$ if and only if $`N`$ is a square mod $`p`$. More precisely, the structure of $`๐ฒ(N,p)`$ is as follows:
* In the cyclotomic case, there is a unique involution $`w`$ on $`X(N,p)`$ such that
$$๐ฒ(N,p)=๐ข(N,p)\times w.$$
* In the non-cyclotomic case,
$$๐ฒ(N,p)\mathrm{PGL}_2(๐ฝ_p).$$
In the first case, the quotient curve $`X(N,p)/w`$ is a Galois covering of $`X^+(N)`$ with group $`๐ข(N,p)`$. In the second case, the quotient of $`X(N,p)`$ by an involution in $`๐ฒ(N,p)`$ is never a Galois covering of $`X^+(N)`$.
*Proof.* Viewed as the quotient $`\mathrm{SL}_2(๐ฝ_p)/\{\pm 1\}`$, the group $`\mathrm{PSL}_2(๐ฝ_p)`$ is generated by the matrices
$$T=(\begin{array}{cc}1& 1\\ 0& 1\end{array}),U=(\begin{array}{cc}1& 0\\ 1& 1\end{array}).$$
On the other hand, the determinant $`\mathrm{GL}_2(๐ฝ_p)๐ฝ_p^{}`$ induces an exact sequence
$$1\mathrm{PSL}_2(๐ฝ_p)\mathrm{PGL}_2(๐ฝ_p)\stackrel{det}{}๐ฝ_p^{}/๐ฝ_{p}^{}{}_{}{}^{2}1,$$
so that $`\mathrm{PSL}_2(๐ฝ_p)`$ can be identified with a subgroup of $`\mathrm{PGL}_2(๐ฝ_p)`$ of index two whose complementary subgroups are those generated by a conjugate of the matrix
$$V=(\begin{array}{cc}0& v\\ 1& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}\end{array}),$$
where $`v`$ is a non-square in $`๐ฝ_p^{}`$. Since one has the relations $`VT=U^{v^1}V`$ and $`VU=T^vV`$โ, a system of generators for $`\mathrm{PGL}_2(๐ฝ_p)`$ is given by $`V`$ and either $`T`$ or $`U`$โ. Now, $`๐ข(N,p)`$ is generated by the automorphisms defined by the matrices
$$T_N=(\begin{array}{cc}1& 1\\ 0& 1\end{array}),U_N=\left(\begin{array}{cc}1& 0\\ \stackrel{~}{N}N& 1\end{array}\right)$$
in $`\mathrm{\Gamma }_0(N)`$, where $`\stackrel{~}{N}`$ is any inverse of $`N`$ mod $`p`$. To give a complementary subgroup for $`๐ข(N,p)`$ inside $`๐ฒ(N,p)`$, let us consider separately the two possibilities for $`N`$ mod $`p`$ :
* If $`N`$ is a square mod $`p`$, then it is also a square mod $`p^2`$. Let $`a,b`$ be any integers satisfying $`a^2Nbp^2=1`$. Then, the matrix
$$Z_N=\left(\begin{array}{cc}aN& bp\\ pN& aN\end{array}\right)$$
defines an involution $`w`$ on $`X(N,p)`$ extending $`w_N`$. Moreover, $`w`$ commutes with the automorphisms defined by $`T_N`$ and $`U_N`$, so it generates a direct cofactor of $`๐ข(N,p)`$ inside $`๐ฒ(N,p)`$. The uniqueness of $`w`$ comes from the fact that $`\mathrm{PSL}_2(๐ฝ_p)`$ has trivial center.
* If $`N`$ is not a square mod $`p`$, then neither is $`\stackrel{~}{N}`$โ. Moreover, the matrix
$$V_N=(\begin{array}{cc}0& 1\\ N& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}\end{array}),$$
which defines an involution on $`X(N,p)`$ extending $`w_N`$, satisfies the relations $`V_NT_N=U_N^NV_N`$ and $`V_NU_N=T_N^{\stackrel{~}{N}}V_N`$ inside $`๐ฒ(N,p)`$. These are precisely the relations that the matrix $`V`$โ, for $`v`$ equal to $`\stackrel{~}{N}`$ mod $`p`$, satisfies with the generators $`T,U`$ of $`\mathrm{PSL}_2(๐ฝ_p)`$. Hence, the group $`๐ฒ(N,p)`$ is isomorphic to $`\mathrm{PGL}_2(๐ฝ_p)`$.
The last assertion in the statement follows from the group structure of $`๐ฒ(N,p)`$ : in the first case, the subgroup $`w`$ is normal, while in the second case $`๐ฒ(N,p)`$ has no normal subgroups of order two because it has trivial center. $`\mathrm{}`$
###### Remark 3.2
The matrices $`Z_N`$ and $`V_N`$ in the proof of Proposition 3.1 have determinant $`N`$โ. Thus, in the same way as the automorphisms in $`๐ข(N,p)`$ are defined by matrices in $`\mathrm{\Gamma }_0(N)`$ acting on $``$, the automorphisms on $`X(N,p)`$ extending $`w_N`$ are defined by matrices in $`\mathrm{M}_2()`$ with determinant $`N`$ and hence lying in $`\mathrm{GL}_2(๐ฝ_p)`$ when reduced mod $`p`$. So we have a mod $`p`$ reduction map
$$๐ฒ(N,p)\mathrm{PGL}_2(๐ฝ_p)$$
whose restriction to $`๐ข(N,p)`$ is the canonical isomorphism onto $`\mathrm{PSL}_2(๐ฝ_p)`$. In the non-cyclotomic case, this map is the isomorphism $`๐ฒ(N,p)\mathrm{PGL}_2(๐ฝ_p)`$ constructed in the proof of Proposition 3.1. We keep this *canonical* isomorphism throughout the rest of the paper.
###### Remark 3.3
In the non-cyclotomic case, all involutions on $`X(N,p)`$ extending $`w_N`$ are conjugated inside $`๐ฒ(N,p)`$. Hence, their defining matrices in $`\mathrm{M}_2()`$ can be obtained conjugating the matrix $`V_N`$ in the proof of Proposition 3.1 by matrices in $`\mathrm{\Gamma }_0(N)`$. So they can be chosen to be of the form
$$(\begin{array}{cc}aN& b\\ cN& aN\end{array}),$$
where $`a,b,c`$ are integers satisfying $`a^2N+bc=1`$. This fact is used in the proof of Proposition 4.3.
In the cyclotomic case, let us write $`X^+(N,p)`$ for the quotient of $`X(N,p)`$ by the only involution $`w`$ in the center of the group $`๐ฒ(N,p)`$. To conclude this section, we give a formula for the genus of $`X(N,p)`$ and compute the values of $`N`$ and $`p`$ for which the curves $`X(N,p)`$ and $`X^+(N,p)`$ have genus zero or one. In the proof of Proposition 3.4, we recall the description of the cusps of $`X_0(N)`$. We refer to \[Gon91\] for this, as well as for the action of the Atkin-Lehner involutions on the set of cusps. Both things are used in the proof of Proposition 3.7.
###### Proposition 3.4
The genus of the modular curve $`X(N,p)`$ is
$$1+\frac{\psi (N)p(p^21)}{24}\frac{p^21}{4}\underset{0<n|N}{}\phi \left((n,N/n)\right),$$
where $`(a,b)`$, $`\phi (r)`$ and $`\psi (N)`$ are the usual notations for the greatest common divisor of the integers $`a`$ and $`b`$, the order of the group $`(/r)^{}`$ and the index of $`\mathrm{\Gamma }_0(N)/\{\pm 1\}`$ in $`\mathrm{PSL}_2()`$, respectively.
*Proof.* The number of cusps of $`X_0(N)`$ is $`\phi (h_n)`$, where the sum is taken over the positive divisors $`n`$ of $`N`$โ, and $`h_n`$ stands for $`(n,N/n)`$. For every divisor $`n`$, there is exactly one cusp for each integer $`m`$ in a system of representatives in $``$ of the group $`(/h_n)^{}`$โ. We just take the integer $`m=1`$ whenever $`\phi (h_n)=1`$. Any such integer $`m`$ can be chosen prime to $`n`$, and the corresponding cusp is then represented by the rational number $`m/n`$. The ramification degree of this cusp over $`X(1)`$ is $`N/(nh_n)`$. On the other hand, the cusps of $`X(p)`$ have ramification degree $`p`$ over $`X(1)`$. Thus, since $`N`$ is prime to $`p`$, the cusps of $`X(N,p)`$ also have ramification degree $`p`$ over $`X_0(N)`$. Moreover, $`X(p)`$ has no elliptic points, so neither has $`X(N,p)`$. Lastly, the degrees of the coverings $`X(N,p)X_0(N)`$ and $`X_0(N)X(1)`$ are $`p(p^21)/2`$ and $`\psi (N)`$, respectively. Hence, the proposition follows from the Hurwitz formula applied to the map $`X(N,p)X(1)`$. $`\mathrm{}`$
###### Corollary 3.5
The modular curve $`X(N,p)`$ has genus greater than one, except for the genus-zero case $`X(2,3)`$ and the elliptic case $`X(4,3)`$.
*Proof.* Since the genera of $`X(p)`$ and $`X_0(N)`$ are greater than one for $`p>5`$ and $`N>49`$, respectively, one only has to check the values that Proposition 3.4 yields in the remaining cases. $`\mathrm{}`$
###### Lemma 3.6
For an odd prime $`p`$ and an integer $`N>1`$ prime to $`p`$, consider the Atkin-Lehner involution $`w_N`$ on the modular curve $`X_0(pN)`$. The only pairs $`(N,p)`$ for which the quotient curve $`X_0(pN)/w_N`$ has genus zero are $`(2,3)`$, $`(4,3)`$, $`(5,3)`$, $`(8,3)`$, $`(11,3)`$, $`(2,5)`$, $`(4,5)`$ and $`(3,7)`$.
*Proof.* For every integer $`D>71`$, the modular curve $`X_0(D)`$ has positive genus and is neither elliptic nor hyperelliptic \[Ogg74\]. For each odd prime $`p`$ and each integer $`N`$ prime to $`p`$ such that $`pN71`$, one can then use the formulae in \[Klu77\] or the tables \[STN92\] to conclude the lemma. $`\mathrm{}`$
###### Proposition 3.7
The curve $`X^+(N,p)`$ has genus greater than one, except for the genus-zero case $`X^+(4,3)`$.
*Proof.* The involution $`w`$, which is defined by the matrix $`Z_N`$ in the proof of Proposition 3.1, restricts to the Atkin-Lehner involution $`w_N`$ on $`X_0(pN)`$, so it induces a Galois covering $`X^+(N,p)X_0(pN)/w_N`$. On the other hand, the cusps of $`X_0(pN)`$ that ramify on $`X(N,p)`$ are those of the form $`m/n`$ with $`p`$ dividing $`n`$, and the ramification degree is always $`p`$ (cf. the proof of Proposition 3.4). In particular, the Hurwitz formula implies that $`X_0(pN)/w_N`$ has genus zero whenever $`X^+(N,p)`$ has genus less than two. By Lemma 3.6, the only pairs $`(N,p)`$, with $`N`$ prime to $`p`$ and square mod $`p`$, for which $`X_0(pN)/w_N`$ has genus zero are $`(4,3)`$ and $`(4,5)`$. In the first case, the involution $`w`$ fixes the cusp $`1/2`$, so $`X^+(4,3)`$ is a genus-zero quotient of the elliptic curve $`X(4,3)`$. Let us now study the second case, for which we consider the following commutative diagram:
The only ramified points of the covering $`X(4,5)X_0(20)`$ are cusps. Moreover, it can be checked that the points lying above the two cusps $`1/2,1/10`$ fixed by $`w_4`$ are also fixed by the involution $`w`$. Thus, the only ramified cusps of the covering $`X^+(4,5)X_0(20)/w_4`$ are the points above $`1/5`$ and $`1/10`$, all of them with ramification degree $`5`$. Then, the Hurwitz formula shows that there must be ten more ramified points, necessarily with ramification degree $`2`$ and lying above the two non-cuspidal points on $`X_0(20)`$ fixed by $`w_4`$, hence the genus of $`X^+(4,5)`$ is four. Notice that there are no other ramified points because the number of points on $`X_0(20)`$ fixed by $`w_4`$ is exactly four (cf. \[Klu77\] or \[STN92\]). $`\mathrm{}`$
## 4 A rational model for the modular curve $`X(N,p)`$
This section deals with the rationality over $``$ for the curve $`X(N,p)`$ as well as for the automorphism group $`๐ฒ(N,p)`$ introduced in the previous section: we fix a certain rational model for $`X(N,p)`$ that makes the automorphisms in $`๐ฒ(N,p)`$ be defined over $`k_p`$. Recall that $`k_p`$ stands for the only quadratic field inside the $`p`$-th cyclotomic extension of $``$. We denote by $`\zeta _p`$ the root of unity $`e^{2\pi i/p}`$.
Since $`X(N,p)`$ is the fiber product of the modular curves $`X(p)`$ and $`X_0(N)`$ over $`X(1)`$, a rational model for the first curve is determined by fixing rational models for the other three curves. Recall that the function field of $`X(1)`$ is generated over $``$ by the elliptic modular function $`j`$. For $`X_0(N)`$, consider the canonical rational model given by the function field $`(j,j_N)`$, where $`j_N`$ is the modular function defined by $`j_N(z)=j(Nz)`$ for $`z`$ in the complex upper-half plane $``$. As for $`X(p)`$, the rational model that we fix satisfies the following property: its extension to $`k_p`$ gives by specialization over an elliptic curve $`E`$ in $`X(1)(\overline{})`$ the fixed field of the projective mod $`p`$ Galois representation $`\overline{\rho }_E`$ attached to the $`p`$-torsion points of $`E`$. This model for $`X(p)`$ is obtained as the following particular case of a general procedure that follows Section II.3 in \[Lig77\] and Section 2 in \[Maz77\].
Fix a non-square $`v`$ in $`๐ฝ_p^{}`$ and take a matrix $`V`$ in $`\mathrm{GL}_2(๐ฝ_p)`$ of order two in $`\mathrm{PGL}_2(๐ฝ_p)`$ and with $`det(V)=v`$. Without risk of confusion, we often identify the matrix $`V`$โ, up to a sign, with its image in $`\mathrm{PGL}_2(๐ฝ_p)`$. Define $`H_V`$ as the inverse image in $`\mathrm{GL}_2(๐ฝ_p)`$ of the subgroup generated by $`V`$ in $`\mathrm{PGL}_2(๐ฝ_p)`$ :
$$H_V=๐ฝ_p^{}๐ฝ_p^{}V.$$
Up to conjugation, $`H_V`$ is the only subgroup of $`\mathrm{GL}_2(๐ฝ_p)`$ containing the center $`๐ฝ_p^{}`$ and reducing inside $`\mathrm{PGL}_2(๐ฝ_p)`$ to a complementary subgroup of $`\mathrm{PSL}_2(๐ฝ_p)`$. The group $`H_V`$ defines, as shown in the following diagram, a rational model $`X_V(p)`$ whose $``$-isomorphism class does not depend on the choice of the matrix $`V`$:
Here $`(\zeta _p)\left(X(p)\right)`$ stands for the field of modular functions for $`\mathrm{\Gamma }(p)`$ whose Fourier expansions have coefficients in $`(\zeta _p)`$. This is a Galois extension of $`\left(X(1)\right)`$ with group $`\mathrm{GL}_2(๐ฝ_p)/\{\pm 1\}`$ and the function field of $`X_V(p)`$ is then the fixed field by the subgroup $`H_V/\{\pm 1\}`$.
Although we should denote by $`X_V(N,p)`$ the rational model for $`X(N,p)`$ obtained from $`X_V(p)`$, we just write $`X(N,p)`$ for simplicity. Without loss of generality, we always take the above non-square $`v`$ equal to $`N^1`$ mod $`p`$ in the non-cyclotomic case. Note that the map $`X(N,p)X_0(N)`$ is defined over $``$ and that the function field $`k_p\left(X(N,p)\right)`$ is a Galois extension of $`\left(X_0(N)\right)`$ with group $`\mathrm{PGL}_2(๐ฝ_p)`$. In particular, the Galois action on the automorphism group $`๐ข(N,p)`$ factors through $`\mathrm{Gal}(k_p/)`$.
The non-cuspidal complex points on $`X(N,p)`$ are in bijection with the isomorphism classes of triples
$$(E,C,[T_1,T_2]_V),$$
where $`E`$ is a complex elliptic curve, $`C`$ is a cyclic subgroup of $`E()`$ of order $`N`$โ, $`[T_1,T_2]`$ is a basis for $`E[p]`$ and $`[T_1,T_2]_V`$ is the corresponding orbit inside $`E[p]\times E[p]`$ by the action of $`H_V`$. Here $`H_V`$ is viewed as a subgroup of automorphisms of $`E[p]`$ through the isomorphism $`\mathrm{GL}_2(๐ฝ_p)\mathrm{Aut}(E[p])`$ fixed by the basis $`[T_1,T_2]`$, so that
$$[T_1,T_2]_V=\left\{[rT_1,rT_2],[rT_1,rT_2]V\right|r๐ฝ_p^{}\}.$$
Two triples of the form $`(E,C,[T_1,T_2]_V)`$ are isomorphic if there is an isomorphism between the corresponding elliptic curves interchanging the cyclic subgroups and the $`H_V`$โorbits.
This bijection is compatible with the usual Galois actions. Thus, a point on $`X(N,p)`$ given by a triple as above with $`j_E0,1728`$ is defined over a number field $`L`$ if and only if the elliptic curve $`E`$ is defined over $`L`$, the subgroup $`C`$ is $`\mathrm{G}_L`$-invariant and the image of the linear Galois representation
$$\rho _E:\mathrm{G}_L\mathrm{GL}_2(๐ฝ_p)$$
attached to $`E[p]`$ lies inside a conjugate of the subgroup $`H_V`$.
We can always assume that the basis $`[T_1,T_2]`$ in a triple $`(E,C,[T_1,T_2]_V)`$ is, inside the corresponding $`H_V`$โorbit, the only one up to a sign that is sent to $`\zeta _p`$ by the Weil pairing. The Galois action on the non-cuspidal points of $`X(N,p)`$ should then be written accordingly: an automorphism $`\sigma `$ of $``$ takes any such a triple to that given by the elliptic curve $`{}_{}{}^{\sigma }E`$, the subgroup $`{}_{}{}^{\sigma }C`$ and the $`H_V`$โorbit of either the basis $`[r^1{}_{}{}^{\sigma }T_{1}^{},r^1{}_{}{}^{\sigma }T_{2}^{}]`$ or the basis $`[(vr)^1{}_{}{}^{\sigma }T_{1}^{},(vr)^1{}_{}{}^{\sigma }T_{2}^{}]V`$โ, depending on whether $`{}_{}{}^{\sigma }\zeta _{p}^{}=\zeta _p^{r^2}`$ or $`{}_{}{}^{\sigma }\zeta _{p}^{}=\zeta _p^{vr^2}`$ for some $`r`$ in $`๐ฝ_p^{}`$, respectively.
The action of the automorphism group $`๐ข(N,p)`$ on the non-cuspidal points of $`X(N,p)`$, and then the Galois action on $`๐ข(N,p)`$, are stated in Proposition 4.1 and Corollary 4.2, respectively. The symbol $`\widehat{}`$ stands henceforth for the matrix (anti)involution given by
$$\widehat{M}=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right){}_{}{}^{\mathrm{t}}M(\begin{array}{cc}0& 1\\ 1& 0\end{array}),$$
where $`{}_{}{}^{\mathrm{t}}M`$ is the transpose of the matrix $`M`$. Alternatively, it can be defined as follows:
$$M=\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)\widehat{M}=(\begin{array}{cc}d& c\\ b& a\end{array}).$$
###### Proposition 4.1
An automorphism in $`๐ข(N,p)`$ represented through the canonical isomorphism $`๐ข(N,p)\mathrm{PSL}_2(๐ฝ_p)`$ by a matrix $`\gamma `$ in $`\mathrm{SL}_2(๐ฝ_p)`$ takes a point $`(E,C,[T_1,T_2]_V)`$ on $`X(N,p)`$ to the point given by the elliptic curve $`E`$, the subgroup $`C`$ and the $`H_V`$โorbit of the $`p`$-torsion basis $`[T_1,T_2]\widehat{\gamma }`$.
*Proof.* Take any matrix $`\left(\genfrac{}{}{0pt}{}{a}{c}\genfrac{}{}{0pt}{}{b}{d}\right)`$ in $`\mathrm{SL}_2()`$ reducing mod $`p`$ to $`\gamma `$. The triple $`(E,C,[T_1,T_2]_V)`$ is isomorphic to one of the form
$$(E_z,1/N,[\mathrm{\hspace{0.17em}1}/p,z/p]_V)$$
for some $`z`$ in $``$, where $`E_z`$ stands for the complex elliptic curve defined by the lattice $`+z`$. The automorphism in the statement sends the pair given by $`z`$ to that given by
$$z^{}=\frac{az+b}{cz+d}.$$
Then, the endomorphism of $``$ defined by multiplication by $`cz+d`$ extends to an isomorphism $`E_z^{}E_z`$ that preserves the subgroup $`1/N`$ and sends the basis $`[\mathrm{\hspace{0.17em}1}/p,z^{}/p]`$ of $`E_z^{}[p]`$ to the basis
$$[(d+cz)/p,(b+az)/p]=[\mathrm{\hspace{0.17em}1}/p,z/p]\widehat{\gamma }$$
of $`E_z[p]`$, so the result follows. $`\mathrm{}`$
###### Corollary 4.2
An automorphism in $`๐ข(N,p)`$ represented through the canonical isomorphism $`๐ข(N,p)\mathrm{PSL}_2(๐ฝ_p)`$ by a matrix $`\gamma `$ in $`\mathrm{SL}_2(๐ฝ_p)`$ is sent by the non-trivial element in $`\mathrm{Gal}(k_p/)`$ to the automorphism in $`๐ข(N,p)`$ corresponding to the matrix $`\widehat{V}\gamma \widehat{V}`$ in $`\mathrm{PSL}_2(๐ฝ_p)`$.
*Proof.* Denote by $`g`$ the automorphism represented by the matrix $`\gamma `$. Take any element $`\sigma `$ in $`\mathrm{G}_{}`$ such that $`{}_{}{}^{\sigma }\zeta _{p}^{}=\zeta _p^v`$ and let $`\gamma _\sigma `$ be a matrix in $`\mathrm{SL}_2(๐ฝ_p)`$ representing the automorphism $`{}_{}{}^{\sigma }g๐ข(N,p)`$. Take also any point $`P`$ in $`X(N,p)(\overline{})`$ given by a triple $`(E,C,[T_1,T_2]_V)`$ with $`j_E0,1728`$. The identity $`{}_{}{}^{\sigma }\left(g(P)\right)={}_{}{}^{\sigma }g(^\sigma P)`$ leads, by means of Proposition 4.1, to an automorphism of the elliptic curve $`{}_{}{}^{\sigma }E`$ interchanging the $`H_V`$โorbits of the $`p`$-torsion bases $`[v^1{}_{}{}^{\sigma }T_{1}^{},v^1{}_{}{}^{\sigma }T_{2}^{}]\widehat{\gamma }V`$ and $`[v^1{}_{}{}^{\sigma }T_{1}^{},v^1{}_{}{}^{\sigma }T_{2}^{}]V\widehat{\gamma }_\sigma `$. This yields the identity $`\widehat{\gamma }_\sigma =V\widehat{\gamma }V`$ in $`\mathrm{PSL}_2(๐ฝ_p)`$. $`\mathrm{}`$
From now on, we fix as follows an involution $`w`$ on $`X(N,p)`$ extending the Atkin-Lehner involution $`w_N`$ on $`X_0(N)`$. Recall that $`๐ฒ(N,p)`$ stands for the group of the Galois covering $`X(N,p)X^+(N)`$, where $`X^+(N)`$ is the quotient of $`X_0(N)`$ by $`w_N`$. In the cyclotomic case, we take as $`w`$ the only involution in the center of $`๐ฒ(N,p)`$ (cf. Proposition 3.1) and denote by $`\sqrt{N}`$ a square root of $`N`$ mod $`p`$. In the non-cyclotomic case, we take as $`w`$ the involution corresponding to the matrix $`\widehat{V}`$ through the canonical isomorphism $`๐ฒ(N,p)\mathrm{PGL}_2(๐ฝ_p)`$ (cf. Remark 3.2). Recall that, in the second case, $`det(V)`$ is taken to be $`N^1`$ mod $`p`$.
###### Proposition 4.3
The involution $`w`$ takes a point $`(E,C,[T_1,T_2]_V)`$ on $`X(N,p)`$ to the point given by the elliptic curve $`E/C`$, the subgroup $`E[N]/C`$ and the $`H_V`$โorbit of the image in $`E/C`$ of the following $`p`$-torsion basis:
* $`[\sqrt{N}^1T_1,\sqrt{N}^1T_2]`$ in the cyclotomic case.
* $`[T_1,T_2]V`$ in the non-cyclotomic case.
*Proof.* According to Remark 3.3 and the proof of Proposition 3.1, the involution $`w`$ is always defined by the action on $``$ of a matrix in $`\mathrm{M}_2()`$ of the form
$$(\begin{array}{cc}aN& b\\ cN& dN\end{array}),$$
with $`adNbc=1`$. Denote by $`\gamma `$ the reduction mod $`p`$ of this matrix. We now proceed as in the proof of Proposition 4.1 : the given triple is isomorphic to one of the form
$$(E_z,1/N,[\mathrm{\hspace{0.17em}1}/p,z/p]_V)$$
for some $`z`$ in $``$, where $`E_z`$ stands for the complex elliptic curve defined by the lattice $`+z`$. The involution $`w`$ sends the triple given by $`z`$ to that given by
$$z^{}=\frac{aNz+b}{cNz+dN}.$$
Then, the endomorphism of $``$ defined by multiplication by $`cz+d`$ extends to an isomorphism
$$E_z^{}E_z/1/N.$$
This isomorphism sends the subgroup $`1/N`$ of $`E_z^{}`$ to the image of $`E_z[N]`$ under the isogeny $`E_zE_z/1/N`$. Also, it sends the basis $`[\mathrm{\hspace{0.17em}1}/p,z^{}/p]`$ of $`E_z^{}[p]`$ to the image of the basis
$$[(d+cz)/p,(N^1b+az)/p]=[\mathrm{\hspace{0.17em}1}/p,z/p]N^1\widehat{\gamma }$$
of $`E_z[p]`$. In the cyclotomic case, $`d=a`$, $`c=p`$ and $`b`$ is a multiple of $`p`$, so that $`a^2`$ equals $`N^1`$ mod $`p`$ and the matrix $`\sqrt{N}^1\widehat{\gamma }`$ is trivial in $`\mathrm{PSL}_2(๐ฝ_p)`$. In the non-cyclotomic case, we have $`\gamma =\pm N\widehat{V}`$. This completes the proof. $`\mathrm{}`$
###### Corollary 4.4
The involution $`w`$ is defined over $``$.
*Proof.* Take any automorphism $`\sigma `$ in $`\mathrm{G}_{}`$. Since $`w_N`$ is defined over $``$, $`{}_{}{}^{\sigma }w`$ is still an involution in $`๐ฒ(N,p)๐ข(N,p)`$. Let $`P`$ be a non-CM point in $`X(N,p)(\overline{})`$ given by a triple $`(E,C,[T_1,T_2]_V)`$. For a fixed model of the elliptic curve $`E/C`$, an isogeny $`\lambda :EE/C`$ with kernel $`C`$ is determined up to a sign. One has the conjugate isogeny $`{}_{}{}^{\sigma }\lambda :{}_{}{}^{\sigma }E{}_{}{}^{\sigma }(E/C)`$. Using Proposition 4.3 and the isomorphism $`{}_{}{}^{\sigma }E/^\sigma C{}_{}{}^{\sigma }(E/C)`$ induced by $`{}_{}{}^{\sigma }\lambda `$, we can verify case by case that $`{}_{}{}^{\sigma }P`$ has the same image by both $`w`$ and $`{}_{}{}^{\sigma }w`$. Consider, for instance, the cyclotomic case and assume $`{}_{}{}^{\sigma }\zeta _{p}^{}=\zeta _p^{r^2}`$ for some $`r`$ in $`๐ฝ_p^{}`$. Then, the point $`{}_{}{}^{\sigma }P`$ is sent to the isomorphism class of the triple given by the elliptic curve $`{}_{}{}^{\sigma }(E/C)`$, the cyclic group $`{}_{}{}^{\sigma }\lambda (^\sigma E[N])`$ and the $`H_V`$โorbit of the basis
$$[(r\sqrt{N})^1{}_{}{}^{\sigma }\lambda (^\sigma T_1),(r\sqrt{N})^1{}_{}{}^{\sigma }\lambda (^\sigma T_2)].$$
By Proposition 4.1, this implies that the matrix in $`\mathrm{PSL}_2(๐ฝ_p)`$ corresponding to the automorphism $`{}_{}{}^{\sigma }ww`$ in $`๐ข(N,p)`$ is the identity, and the result follows. $`\mathrm{}`$
###### Remark 4.5
We can conclude that the Galois covering $`X(N,p)X^+(N)`$ is defined over $`k_p`$. In other words, the function field $`k_p\left(X(N,p)\right)`$ is a Galois extension of $`k_p\left(X^+(N)\right)`$, with group (anti)isomorphic to $`๐ฒ(N,p)`$. As a matter of fact, $`k_p\left(X(N,p)\right)`$ is a Galois extension of $`\left(X^+(N)\right)`$.
Let us finish this section by reviewing the moduli interpretation of the rational points on $`X^+(N)`$. The non-cuspidal points of $`X_0(N)()`$ are in bijection with the isomorphism classes of pairs $`(E,C)`$, where $`E`$ is a complex elliptic curve and $`C`$ is a cyclic subgroup of $`E()`$ of order $`N`$โ. Such a point is defined over a number field $`L`$ if and only if $`E`$ and $`C`$ are defined over $`L`$, which means that $`{}_{}{}^{\sigma }E=E`$ and $`{}_{}{}^{\sigma }C=C`$ for all $`\sigma `$ in $`\mathrm{G}_L`$. A point on $`X(N,p)`$ given by a triple $`(E,C,[T_1,T_2]_V)`$ has image on $`X_0(N)`$ given by the pair $`(E,C)`$. In particular, the involution $`w_N`$ sends this pair to $`(E/C,E[N]/C)`$.
Let $`E`$ be an elliptic curve defined over $`L`$, and let $`EE^{}`$ be an isogeny with cyclic kernel $`C`$ of order $`N`$โ. Assume that $`E`$ has no CM, so that an isogeny from $`E`$ to any elliptic curve is determined up to a sign by its degree. Then, the subgroup $`C`$ is defined over $`L`$ if and only if $`E^{}`$ admits a model over $`L`$.
Now, suppose that $`E`$ and $`E^{}`$ are defined over a quadratic field $`k`$, so that the pair $`(E,C)`$ defines a $`k`$-rational point $`P`$ on $`X_0(N)`$. This point is rational if and only if both $`E`$ and $`E^{}`$ have a model over $``$. In this case we say that the couple $`\{E,E^{}\}`$ is a *fake $``$-curve of degree $`N`$*. Otherwise, the image of $`P`$ on $`X^+(N)`$ is rational if and only if $`E^{}`$ is isomorphic to the Galois conjugate $`{}_{}{}^{\nu }E`$ of $`E`$. Indeed, since $`E`$ has no CM, an isogeny $`\mu :E{}_{}{}^{\nu }E`$ with kernel $`C`$ sends $`E[N]`$ to $`{}_{}{}^{\nu }C`$, so the existence of such an isogeny $`\mu `$ amounts to the equality $`w_N(P)={}_{}{}^{\nu }P`$ in $`X_0(N)(k)`$. Thus, every non-cuspidal non-CM rational point on $`X^+(N)`$ comes from a pair $`(E,C)`$ on $`X_0(N)`$ defined over some quadratic field and yielding a (possibly fake) $``$-curve of degree $`N`$โ.
## 5 The twisted curves in the cyclotomic case
Assume $`N`$ to be a square mod $`p`$. The structure of this section is as follows. We first obtain from a modular point of view the fixed field of the Galois representation $`\varrho _E`$ attached in Section 2 to a $``$-curve $`E`$ of degree $`N`$โ. Next, we produce the twisted modular curves whose non-cuspidal non-CM rational points give the $``$-curves of degree $`N`$ realizing a fixed projective mod $`p`$ Galois representation with cyclotomic determinant. We also include a result on the finiteness of the number of such $``$-curves.
Recall that $`X^+(N,p)`$ denotes the quotient of $`X(N,p)`$ by the involution $`w`$. The induced map $`X^+(N,p)X^+(N)`$ is a Galois covering with automorphism group $`๐ข(N,p)`$ and hence defined over $`k_p`$ (cf. Proposition 3.1 and Remark 4.5). The function field $`k_p\left(X^+(N,p)\right)`$ is in fact a Galois extension of $`\left(X^+(N)\right)`$, with group $`\mathrm{PGL}_2(๐ฝ_p)`$ :
###### Proposition 5.1
The function field $`k_p\left(X^+(N,p)\right)`$ produces, by specialization over a rational point on $`X^+(N)`$ corresponding to a $``$-curve $`E`$, the fixed field of the Galois representation $`\varrho _E`$.
*Proof.* Let $`E`$ be a $``$-curve of degree $`N`$ defined over a quadratic field $`k`$. Fix an automorphism $`\nu `$ in $`\mathrm{G}_{}\mathrm{G}_k`$ and an isogeny $`\mu :E{}_{}{}^{\nu }E`$ of degree $`N`$โ. If we let $`C`$ be the kernel of $`\mu `$, the pair $`(E,C)`$ defines a $`k`$-rational point on $`X_0(N)`$ with rational image on $`X^+(N)`$. The preimages on $`X^+(N,p)`$ of this rational point are given by the couples $`\{P,w(P)\}`$ for all points $`P`$ on $`X(N,p)`$ represented by a triple of the form $`(E,C,[T_1,T_2]_V)`$. If we denote by $`H`$ the subgroup of $`\mathrm{G}_{k_p}`$ fixing those couples, what the proposition asserts is that $`H`$ equals the kernel of $`\varrho _E`$. This kernel is indeed a subgroup of $`\mathrm{G}_{k_p}`$ because the fixed field of $`det\varrho _E`$ is $`k_p`$ (cf. Corollary 2.5). For a point $`P`$ as above, $`w(P)`$ is given by the triple $`({}_{}{}^{\nu }E,{}_{}{}^{\nu }C,[\sqrt{N}^1\mu (T_1),\sqrt{N}^1\mu (T_2)]_V)`$. Take now any $`\sigma `$ in $`\mathrm{G}_{k_p}`$, so that $`{}_{}{}^{\sigma }\zeta _{p}^{}=\zeta _p^{r^2}`$ for some $`r`$ in $`๐ฝ_p^{}`$. If $`\sigma \mathrm{G}_k`$, then $`\sigma H`$ if and only if $`{}_{}{}^{\sigma }P=P`$โ, namely if and only if $`{}_{}{}^{\sigma }T=\pm rT`$ for all points $`T`$ in $`E[p]`$. If $`\sigma \mathrm{G}_k`$, then $`\sigma H`$ if and only if $`{}_{}{}^{\sigma }P=w(P)`$, namely if and only if $`{}_{}{}^{\sigma }T=\pm r\sqrt{N}^1\mu (T)`$ for all points $`T`$ in $`E[p]`$. Therefore, the result follows from the definition of $`\varrho _E`$. $`\mathrm{}`$
Suppose that we are now given a Galois representation
$$\varrho :\mathrm{G}_{}\mathrm{PGL}_2(๐ฝ_p)$$
with cyclotomic determinant, which means that the fixed field of $`det\varrho `$ is $`k_p`$. For the moduli problem of classifying the $``$-curves of degree $`N`$ realizing $`\varrho `$, we twist the curve $`X^+(N,p)`$ by certain elements in the cohomology set $`H^1(\mathrm{G}_{},๐ข(N,p))`$. Recall that the twists of a curve defined over $``$, up to $``$-isomorphism, are in bijection with the elements in the first cohomology set of $`\mathrm{G}_{}`$ with values in the automorphism group of the curve.
The Galois action on $`๐ข(N,p)`$ is known from Corollary 4.2. Now, the action by conjugation of $`\mathrm{PGL}_2(๐ฝ_p)`$ makes this group isomorphic to the automorphism group of $`\mathrm{PSL}_2(๐ฝ_p)`$. Hence, the canonical isomorphism $`๐ข(N,p)\mathrm{PSL}_2(๐ฝ_p)`$ induces an isomorphism $`\mathrm{Aut}\left(๐ข(N,p)\right)\mathrm{PGL}_2(๐ฝ_p)`$ through which the Galois action on $`๐ข(N,p)`$ can be described by the morphism
$$\eta :\mathrm{G}_{}\mathrm{Gal}(k_p/)\widehat{V}\mathrm{PGL}_2(๐ฝ_p).$$
Consider then the $`1`$-cocycles $`\xi =\varrho _{}\eta `$ and $`\xi ^{}=\varrho _{}^{}\eta `$, where $`\varrho _{}(\sigma )={}_{}{}^{\mathrm{t}}\varrho (\sigma ^1)`$ and $`\varrho _{}^{}(\sigma )=\widehat{V}\varrho _{}(\sigma )\widehat{V}`$ for all $`\sigma `$ in $`\mathrm{G}_{}`$. The cyclotomic hypothesis allows us to regard them, through the above canonical isomorphism, as cocycles with values in $`๐ข(N,p)`$. The cocycle condition for $`\xi `$, namely $`\xi _{\sigma \tau }=\xi _\sigma {}_{}{}^{\sigma }\xi _{\tau }^{}`$ for all $`\sigma ,\tau `$ in $`\mathrm{G}_{}`$, can be easily checked case by case, depending on whether $`\sigma `$ and $`\tau `$ belong to $`\mathrm{G}_{k_p}`$ or not. The same holds for $`\xi ^{}`$โ. The cocycle $`\xi `$ defines a rational model $`X^+(N,p)_\varrho `$ for the corresponding twist of $`X^+(N,p)`$, together with an isomorphism
$$\psi _+:X^+(N,p)_\varrho X^+(N,p)$$
satisfying $`\psi _+=\xi _\sigma {}_{}{}^{\sigma }\psi _{+}^{}`$ for every $`\sigma `$ in $`\mathrm{G}_{}`$. Let us denote by $`X^+(N,p)_\varrho ^{}`$ and $`\psi _+^{}`$ the analogous twist and isomorphism defined by the cocycle $`\xi ^{}`$โ.
###### Theorem 5.2
There exists a (possibly fake) $``$-curve of degree $`N`$ realizing $`\varrho `$ if and only if the set of non-cuspidal non-CM rational points on the curves $`X^+(N,p)_\varrho `$ and $`X^+(N,p)_\varrho ^{}`$ is not empty. In this case, the compositions of the isomorphisms $`\psi _+`$ and $`\psi _+^{}`$ with the natural map $`X^+(N,p)X^+(N)`$ define a surjective map from this set of points to the set of isomorphism classes of :
$``$$``$-curves of degree $`N`$ up to Galois conjugation realizing $`\varrho `$,
$``$ fake $``$-curves of degree $`N`$ realizing $`\varrho `$.
This map is bijective if and only if the centralizer in $`\mathrm{PGL}_2(๐ฝ_p)`$ of the image of $`\varrho `$ is trivial.
*Proof.* The rational points on $`X^+(N,p)_\varrho `$ correspond via $`\psi _+`$ to the couples of the form $`\{P,w(P)\}`$, where $`P`$ is an algebraic point on $`X(N,p)`$ such that, for each given automorphism $`\sigma `$ in $`\mathrm{G}_{}`$, either $`\xi _\sigma ({}_{}{}^{\sigma }P)=P`$ or $`\xi _\sigma ({}_{}{}^{\sigma }P)=w(P)`$.
Let $`P`$ be a non-CM point in $`X(N,p)(\overline{})`$ given by a triple $`(E,C,[T_1,T_2]_V)`$ . We use the basis $`[T_1,T_2]`$ of $`E[p]`$ to fix the isomorphism $`\mathrm{Aut}\left(E[p]\right)\mathrm{GL}_2(๐ฝ_p)`$. By virtue of Proposition 4.1, the condition $`{}_{}{}^{\sigma }P=\xi _\sigma ^1(P)`$ for all $`\sigma `$ in $`\mathrm{G}_{}`$ amounts to saying that $`\{E,E/C\}`$ is a fake $``$-curve of degree $`N`$ such that the equality
$$\varrho _E(\sigma )=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\varrho (\sigma )\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)$$
(1)
holds in $`\mathrm{PGL}_2(๐ฝ_p)`$ for every $`\sigma `$ in $`\mathrm{G}_{}`$. Here, we extend the notation $`\varrho _E`$ to the case of elliptic curves over $``$ by putting $`\varrho _E=\overline{\rho }_E`$.
If, on the other hand, there exists $`\nu `$ in $`\mathrm{G}_{}`$ for which $`\xi _\nu ({}_{}{}^{\nu }P)=w(P)`$, then $`E`$ must be a quadratic $``$-curve for the point $`\psi _+^1\left(\{P,w(P)\}\right)`$ on $`X^+(N,p)_\varrho `$ to be rational. Indeed, in this case the subgroup of $`\mathrm{G}_{}`$ consisting of those automorphisms $`\sigma `$ which satisfy $`\xi _\sigma ({}_{}{}^{\sigma }P)=P`$ has index two, so it is of the form $`\mathrm{G}_k`$ for some quadratic field $`k`$, and then the condition $`{}_{}{}^{\sigma }P=\xi _\sigma ^1(P)`$ for all $`\sigma `$ in $`\mathrm{G}_k`$ forces the elliptic curve $`E`$ and the subgroup $`C`$ to be defined over $`k`$, while the condition $`w({}_{}{}^{\nu }P)=\xi _\nu ^1(P)`$ gives an isogeny $`\lambda :{}_{}{}^{\nu }EE`$ with kernel $`{}_{}{}^{\nu }C`$.
So assume now $`E`$ and $`C`$ to be defined over a quadratic field $`k`$ and let $`\lambda `$ be an isogeny as above. Then, for $`\sigma \mathrm{G}_k`$, the point $`w(^\sigma P)`$ is represented by the triple given by the elliptic curve $`E`$, the cyclic group $`C`$ and the $`H_V`$โorbit of the basis
This comes from Proposition 4.3 and the isomorphism $`{}_{}{}^{\nu }E/{}_{}{}^{\nu }CE`$ induced by the isogeny $`\lambda `$. Notice that the second case does not occur whenever $`k=k_p`$.
On the other hand, the automorphism $`\xi _\sigma ^1`$ is given by $`{}_{}{}^{\mathrm{t}}\varrho (\sigma )`$, if $`\sigma \mathrm{G}_{k_p}`$, or by $`\widehat{V}{}_{}{}^{\mathrm{t}}\varrho (\sigma )`$, if $`\sigma \mathrm{G}_{k_p}`$. Then, by applying Proposition 4.1 to each case, we obtain that the point $`\psi _+^1\left(\{P,w(P)\}\right)`$ on $`X^+(N,p)_\varrho `$ is rational if and only if condition (1) holds for every $`\sigma `$ in $`\mathrm{G}_{}`$.
Similarly, consider a point on $`X^+(N,p)_\varrho ^{}`$ corresponding via $`\psi _+^{}`$ to a point on $`X^+(N,p)`$ obtained from a triple $`(E,C,[T_1,T_2]_V)`$. By the same reasoning as above, this point is rational if and only if the pair $`(E,C)`$ represents a (possibly fake) $``$-curve of degree $`N`$ such that, for every $`\sigma `$ in $`\mathrm{G}_{}`$,
$$\varrho _E(\sigma )=V\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\varrho (\sigma )\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)V.$$
(2)
Let us now consider a (possibly fake) $``$-curve of degree $`N`$ given by some point $`(E,C)`$ on $`X_0(N)`$ and assume $`\varrho _E=\varrho `$. Since this is an equality up to conjugation in $`\mathrm{PGL}_2(๐ฝ_p)`$, it amounts to the existence of a basis $`[T_1,T_2]`$ of $`E[p]`$ for which condition (1) holds for every $`\sigma `$ in $`\mathrm{G}_{}`$. Moreover, we can suppose that such a basis is sent to either $`\zeta _p`$ or $`\zeta _p^{v^1}`$ by the Weil pairing. In the first case, the image on $`X^+(N,p)`$ of the triple $`(E,C,[T_1,T_2]_V)`$ defines through $`\psi _+`$ a rational point on $`X^+(N,p)_\varrho `$. In the second case, take $`[T_1^{},T_2^{}]=[T_1,T_2]V`$. For this new basis, which is sent to $`\zeta _p`$ under the Weil pairing, condition (2) is satisfied for every $`\sigma `$ in $`\mathrm{G}_{}`$. So the image on $`X^+(N,p)`$ of the triple $`(E,C,[T_1^{},T_2^{}]_V)`$ defines through $`\psi _+^{}`$ a rational point on $`X^+(N,p)_\varrho ^{}`$.
This proves the first part of the statement, including the surjectivity of the map whenever it is defined. To discuss its injectivity, consider a point $`(E,C)`$ on $`X_0(N)`$ yielding a (possibly fake) $``$-curve of degree $`N`$โ. Suppose that one can take two different rational points on the twists, corresponding (via $`\psi _+`$ or $`\psi _+^{}`$) to points on $`X^+(N,p)`$ obtained from two triples of the form $`(E,C,[T_1,T_2]_V)`$. Three different cases must be distinguished to complete the proof:
* Both rational points are on $`X^+(N,p)_\varrho `$ if and only if there is a non-trivial element $`\gamma `$ in $`\mathrm{PSL}_2(๐ฝ_p)`$, representing a basis change in $`E[p]`$, such that
$$\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\varrho (\sigma )\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)=\gamma \left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\varrho (\sigma )\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\gamma ^1$$
for all $`\sigma `$ in $`\mathrm{G}_{}`$. This amounts to the existence of a non-trivial element in $`\mathrm{PSL}_2(๐ฝ_p)`$ commuting with all the elements in the image of $`\varrho `$.
* The same characterization is obtained if both points lie on $`X^+(N,p)_\varrho ^{}`$.
* One of the points is on $`X^+(N,p)_\varrho `$ and the other on $`X^+(N,p)_\varrho ^{}`$ if and only if there exists $`\gamma `$ in $`\mathrm{PSL}_2(๐ฝ_p)`$ such that
$$V\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\varrho (\sigma )\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)V=\gamma \left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\varrho (\sigma )\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\gamma ^1$$
for every $`\sigma `$ in $`\mathrm{G}_{}`$. This amounts to the existence of an element in $`\mathrm{PGL}_2(๐ฝ_p)`$ not lying in $`\mathrm{PSL}_2(๐ฝ_p)`$ and commuting with all the elements in the image of $`\varrho `$. $`\mathrm{}`$
Note that the set of points in Theorem 5.2 is always finite whenever the genus of $`X^+(N)`$ is greater than one. One can assure this for $`N>131`$ : indeed, the modular curve $`X_0(N)`$ has genus at least two and is neither hyperelliptic \[Ogg74\] nor bielliptic \[Bar99\] for any such integer $`N`$โ. Using Proposition 3.7, one actually gets the following improvement.
###### Corollary 5.3
For $`N`$ square mod $`p`$, the number of isomorphism classes of $``$-curves of degree $`N`$ realizing $`\varrho `$ is finite, except possibly in the case $`N=4`$, $`p=3`$.
For $``$-curves of degree $`N`$ realizing $`\varrho `$, a different moduli description that gets rid of fake $``$-curves can be given for every quadratic field of definition. In order to do that, we twist $`X(N,p)`$ by two certain elements in the cohomology set $`H^1(\mathrm{G}_{},๐ฒ(N,p))`$ that are naturally obtained from the above cocycles $`\xi `$ and $`\xi ^{}`$ as follows. By Proposition 3.1 and Corollary 4.4, the $`\mathrm{G}_{}`$-group $`๐ฒ(N,p)`$ equals the direct product of $`\mathrm{G}_{}`$-groups $`๐ข(N,p)\times w`$. Then, $`H^1(\mathrm{G}_{},๐ฒ(N,p))`$ is also the direct product of the corresponding cohomology sets. Fix now a quadratic field $`k`$ and take the Galois character
$$\chi _k:\mathrm{G}_{}\mathrm{Gal}(k/)w.$$
We then consider the $`1`$-cocycle $`\xi \chi _k`$ and the rational model $`X(N,p)_{\varrho ,k}`$ for the corresponding twist, along with the isomorphism
$$\psi _k:X(N,p)_{\varrho ,k}X(N,p)$$
satisfying $`\psi _k=(\xi \chi _k)_\sigma {}_{}{}^{\sigma }\psi _{k}^{}`$ for every $`\sigma `$ in $`\mathrm{G}_{}`$. Analogously, let us denote by $`X(N,p)_{\varrho ,k}^{}`$ and $`\psi _k^{}`$ the twist and the isomorphism defined by the cocycle $`\xi ^{}\chi _k`$ .
###### Theorem 5.4
There exists a $``$-curve of degree $`N`$ defined over $`k`$ realizing $`\varrho `$ if and only if the set of non-cuspidal non-CM rational points on the curves $`X(N,p)_{\varrho ,k}`$ and $`X(N,p)_{\varrho ,k}^{}`$ is not empty. In this case, the compositions of the isomorphisms $`\psi _k`$ and $`\psi _k^{}`$ with the natural map $`X(N,p)X_0(N)`$ define a surjective map from this set of points to the set of isomorphism classes of $``$-curves of degree $`N`$ defined over $`k`$ realizing $`\varrho `$. This map is bijective if and only if the centralizer in $`\mathrm{PGL}_2(๐ฝ_p)`$ of the image of $`\varrho `$ is trivial.
*Proof.* The rational points on $`X_{\varrho ,k}(N,p)`$ correspond via $`\psi _k`$ to the algebraic points $`P`$ on $`X(N,p)`$ such that
$$\xi _\sigma ^1(P)=\{\begin{array}{cc}{}_{}{}^{\sigma }P\hfill & \mathrm{for}\sigma \mathrm{G}_k,\hfill \\ w(^\sigma P)\hfill & \mathrm{for}\sigma \mathrm{G}_k.\hfill \end{array}$$
The proof runs then in a very similar way to that of Theorem 5.2, so we omit the details. In the current case, a non-CM point on $`X(N,p)_{\varrho ,k}`$ corresponding via $`\psi _k`$ to a triple $`(E,C,[T_1,T_2]_V)`$ is rational if and only if $`E`$ is defined over $`k`$, there exists an isogeny from $`E`$ to its Galois conjugate with kernel $`C`$ and condition (1) holds for every $`\sigma `$ in $`\mathrm{G}_{}`$ whenever one uses the basis $`[T_1,T_2]`$ to fix the isomorphism $`\mathrm{Aut}\left(E[p]\right)\mathrm{GL}_2(๐ฝ_p)`$. The same characterization is valid for the rational points on $`X(N,p)_{\varrho ,k}^{}`$ if we replace $`\psi _k`$ by $`\psi _k^{}`$ and condition (1) by condition (2). $`\mathrm{}`$
###### Remark 5.5
One can check that $`\xi `$ and $`\xi ^{}`$ are cohomologous as $`1`$-cocycles with values in $`๐ข(N,p)`$ if and only if the centralizer in $`\mathrm{PGL}_2(๐ฝ_p)`$ of the image of $`\varrho `$ does not lie in $`\mathrm{PSL}_2(๐ฝ_p)`$. Thus, the twists $`X^+(N,p)_\varrho `$ and $`X^+(N,p)_\varrho ^{}`$ are not a priori isomorphic over $``$. The same holds for the twisted curves $`X(N,p)_{\varrho ,k}`$ and $`X(N,p)_{\varrho ,k}^{}`$ . Moreover, it can be shown that the involution $`w`$ does not switch the rational points on $`X(N,p)_{\varrho ,k}`$ and $`X(N,p)_{\varrho ,k}^{}`$ , so finding the underlying $``$-curves requires in general the rational points on both twists.
## 6 The twisted curve in the non-cyclotomic case
Assume $`N`$ to be a non-square mod $`p`$. This section is the analogue of the previous one for the non-cyclotomic case. Unlike the cyclotomic case, now the quadratic field of definition for the potential $``$-curves of degree $`N`$ realizing a given projective mod $`p`$ Galois representation is fixed by the determinant. Moreover, only one twist is needed for the moduli classification of such $``$-curves. We prove this in Theorem 6.4 below. For the sake of completeness, let us begin as before with the modular construction of the fixed field of the Galois representation $`\varrho _E`$ attached to a $``$-curve $`E`$ of degree $`N`$โ. The procedure is now more intricate.
Recall that the group $`๐ฒ(N,p)`$ of the covering $`X(N,p)X^+(N)`$ is canonically isomorphic to $`\mathrm{PGL}_2(๐ฝ_p)`$. The action by conjugation of this group makes it isomorphic to its automorphism group. Thus, by virtue of Corollary 4.2 and Corollary 4.4, the Galois action on $`๐ฒ(N,p)`$ is given by the morphism
$$\eta :\mathrm{G}_{}\mathrm{Gal}(k_p/)w๐ฒ(N,p),$$
where we identify $`๐ฒ(N,p)`$ with its (inner) automorphism group.
Let $`\stackrel{~}{X}(N,p)`$ be the twist of $`X(N,p)`$ defined by the $`1`$-cocycle $`\eta `$. Likewise, denote by $`\stackrel{~}{X}_0(N)`$ the twist of $`X_0(N)`$ defined by the $`1`$-cocycle
$$\mathrm{G}_{}\mathrm{Gal}(k_p/)w_N.$$
We write $`\stackrel{~}{X}^+(N)`$ for the quotient of $`\stackrel{~}{X}_0(N)`$ by the involution corresponding to $`w_N`$. Consider the following commutative diagram, where the morphisms are the natural ones:
As remarked in the proof of the next lemma, the isomorphism $`\stackrel{~}{X}^+(N)X^+(N)`$ is actually defined over $``$.
###### Lemma 6.1
The Galois covering $`\stackrel{~}{X}(N,p)\stackrel{~}{X}^+(N)`$ is defined over $``$.
*Proof.* Denote by $`\varphi :\stackrel{~}{X}(N,p)X(N,p)`$ and $`\varphi _0:\stackrel{~}{X}_0(N)X_0(N)`$ the isomorphisms in the above diagram. They are defined over $`k_p`$ and satisfy $`{}_{}{}^{\sigma }\varphi \varphi ^1=w`$ and $`{}_{}{}^{\sigma }\varphi _{0}^{}\varphi _0^1=w_N`$ for $`\sigma \mathrm{G}_{k_p}`$. Then, the involution $`\varphi _0^1w_N\varphi _0`$ on $`\stackrel{~}{X}_0(N)`$ is defined over $``$. Hence, so is the corresponding quotient map $`\stackrel{~}{X}_0(N)\stackrel{~}{X}^+(N)`$. The isomorphism $`\varphi _+:X^+(N)\stackrel{~}{X}^+(N)`$ induced by $`\varphi _0^1`$ sends a couple $`\{P,w_N(P)\}`$ to $`\{\varphi _0^1(P),\varphi _0^1w_N(P)\}`$. It is easily checked to satisfy $`{}_{}{}^{\sigma }\varphi _{+}^{}=\varphi _+`$ for all $`\sigma `$ in $`\mathrm{G}_{}`$. The same is true for the morphism $`\stackrel{~}{X}(N,p)\stackrel{~}{X}_0(N)`$ induced from the natural map $`X(N,p)X_0(N)`$ by the isomorphisms $`\varphi `$ and $`\varphi _0`$. Finally, the automorphisms of the covering $`\stackrel{~}{X}(N,p)\stackrel{~}{X}^+(N)`$ are also defined over $``$. Indeed, the relation $`{}_{}{}^{\sigma }(\varphi ^1\vartheta \varphi )=\varphi ^1w(w\vartheta w)w\varphi =\varphi ^1\vartheta \varphi `$ holds for $`\vartheta ๐ฒ(N,p)`$ and $`\sigma \mathrm{G}_{k_p}`$. For another proof of the existence of such a rational covering we refer to \[Shi74\]. $`\mathrm{}`$
###### Remark 6.2
The function field of $`\stackrel{~}{X}(N,p)`$ over $``$ is identified, through the isomorphism $`\varphi `$ in the proof of Lemma 6.1, with a subfield of $`k_p\left(X(N,p)\right)`$. As shown in the following diagram, it is a Galois extension of $`\left(X^+(N)\right)`$ with group isomorphic to $`\mathrm{PGL}_2(๐ฝ_p)`$ :
###### Proposition 6.3
The function field $`(\stackrel{~}{X}(N,p))`$ gives, by specialization over a rational point on $`X^+(N)`$ corresponding to a $``$-curve $`E`$, the fixed field of the representation $`\varrho _E`$.
*Proof.* With the same notations as in the proof of Proposition 5.1, take a cyclic isogeny $`\mu :E{}_{}{}^{\nu }E`$ with kernel $`C`$ of order $`N`$โ. Consider the isomorphism $`\varphi :\stackrel{~}{X}(N,p)X(N,p)`$ in the proof of Lemma 6.1. Let $`H`$ be the subgroup of $`\mathrm{G}_{}`$ fixing the points on $`\stackrel{~}{X}(N,p)`$ corresponding through $`\varphi `$ to the points on $`X(N,p)`$ of the form $`P`$ or $`w(P)`$, where $`P`$ is given by a triple of the form $`(E,C,[T_1,T_2]_V)`$. We must show that $`H`$ is the kernel of $`\varrho _E`$. For a point $`P`$ as above, $`w(P)`$ is represented by the triple given by $`{}_{}{}^{\nu }E`$, $`{}_{}{}^{\nu }C`$ and the $`H_V`$โorbit of the basis $`[\mu (T_1),\mu (T_2)]V`$โ. Using the definition of $`\varphi `$, we see that the group $`H`$ consists of those $`\sigma \mathrm{G}_{k_p}`$ satisfying $`{}_{}{}^{\sigma }P=P`$ and those $`\sigma \mathrm{G}_{k_p}`$ satisfying $`{}_{}{}^{\sigma }P=w(P)`$. Moreover, any such $`\sigma `$ lies in $`\mathrm{G}_k`$ if and only if it lies in $`\mathrm{G}_{k_p}`$. Take now any automorphism $`\sigma `$ in $`\mathrm{G}_{}`$. If $`{}_{}{}^{\sigma }\zeta _{p}^{}=\zeta _p^{r^2}`$ for some $`r`$ in $`๐ฝ_p^{}`$, then $`\sigma H`$ if and only if $`{}_{}{}^{\sigma }P=P`$โ, namely if and only if $`{}_{}{}^{\sigma }T=\pm rT`$ for all points $`T`$ in $`E[p]`$. If $`{}_{}{}^{\sigma }\zeta _{p}^{}=\zeta _p^{r^2N^1}`$ for some $`r`$ in $`๐ฝ_p^{}`$, then $`\sigma H`$ if and only if $`{}_{}{}^{\sigma }P=w(P)`$, namely if and only if $`{}_{}{}^{\sigma }T=\pm rN^1\mu (T)`$ for all points $`T`$ in $`E[p]`$. So the result follows from the definition of $`\varrho _E`$. $`\mathrm{}`$
Suppose that we have now a Galois representation
$$\varrho :\mathrm{G}_{}\mathrm{PGL}_2(๐ฝ_p)$$
with non-cyclotomic determinant. Recall that any $``$-curves of degree $`N`$ realizing $`\varrho `$ must be defined over the fixed field of $`\epsilon det\varrho `$, where $`\epsilon `$ is the character attached to $`k_p`$ (cf. Corollary 2.5). Denote this quadratic field by $`k`$. For the moduli classification of such $``$-curves, we produce a twist of $`X(N,p)`$ from a certain element in the cohomology set $`H^1(\mathrm{G}_{},๐ฒ(N,p))`$, as follows. The canonical isomorphism $`๐ฒ(N,p)\mathrm{PGL}_2(๐ฝ_p)`$ allows us to regard the projective representation $`\varrho _{}`$ in Section 5 as a morphism taking values in $`๐ฒ(N,p)`$. As before, let $`\eta `$ stand for the morphism giving the Galois action on $`๐ฒ(N,p)`$. Then, consider the $`1`$-cocycle $`\xi =\varrho _{}\eta `$. For the twist of $`X(N,p)`$ defined by $`\xi `$, we fix a rational model $`X(N,p)_\varrho `$ along with an isomorphism
$$\psi :X(N,p)_\varrho X(N,p)$$
satisfying $`\psi =\xi _\sigma {}_{}{}^{\sigma }\psi `$ for every $`\sigma `$ in $`\mathrm{G}_{}`$.
###### Theorem 6.4
There exists a $``$-curve of degree $`N`$ realizing $`\varrho `$ if and only if the set of non-cuspidal non-CM rational points on the curve $`X(N,p)_\varrho `$ is not empty. In this case, the composition of $`\psi `$ with the natural map $`X(N,p)X^+(N)`$ defines a surjective map from this set of points to the set of isomorphism classes of $``$-curves of degree $`N`$ up to Galois conjugation realizing $`\varrho `$. This map is bijective if and only if the centralizer in $`\mathrm{PGL}_2(๐ฝ_p)`$ of the image of $`\varrho `$ is trivial.
*Proof.* The first part of the proof goes along the lines of those of Theorem 5.2 and Theorem 5.4. Let us fix an automorphism $`\nu `$ in $`\mathrm{G}_{}\mathrm{G}_k`$. The rational points on $`X_\varrho (N,p)`$ correspond via $`\psi `$ to the algebraic points $`P`$ on $`X(N,p)`$ satisfying
$$\varrho _{}(\sigma )^1(P)=\{\begin{array}{cc}{}_{}{}^{\sigma }P\hfill & \mathrm{for}\sigma \mathrm{G}_{k_p},\hfill \\ w(^\sigma P)\hfill & \mathrm{for}\sigma \mathrm{G}_{k_p}.\hfill \end{array}$$
Note that the automorphism $`\varrho _{}(\sigma )^1`$ belongs to $`๐ข(N,p)`$ if and only if $`\sigma `$ lies in either both $`\mathrm{G}_k`$ and $`\mathrm{G}_{k_p}`$ or none of them. In particular, a non-CM point $`P`$ given by a triple $`(E,C,[T_1,T_2]_V)`$ can satisfy the above condition only if $`E`$ and $`C`$ are defined over $`k`$ and there is an isogeny $`\lambda :{}_{}{}^{\nu }EE`$ with kernel $`{}_{}{}^{\nu }C`$. With these hypotheses on $`E`$ and $`C`$, and for $`\sigma \mathrm{G}_k`$, the point $`w(^\sigma P)`$ is represented by the triple given by $`E`$, $`C`$ and the $`H_V`$โorbit of the basis
In the second case, and also for $`\sigma \mathrm{G}_k\mathrm{G}_{k_p}`$, the automorphism $`\varrho _{}(\sigma )^1`$ is given by the matrix $`{}_{}{}^{\mathrm{t}}\varrho (\sigma )`$ in $`\mathrm{PSL}_2(๐ฝ_p)`$. In the other case, and also for $`\sigma \mathrm{G}_k\mathrm{G}_{k_p}`$, the automorphism $`w\varrho _{}(\sigma )^1`$ is given by the matrix $`\widehat{V}{}_{}{}^{\mathrm{t}}\varrho (\sigma )`$ in $`\mathrm{PSL}_2(๐ฝ_p)`$. So, taking as in the proof of Theorem 5.2 the basis $`[T_1,T_2]`$ to fix the isomorphism $`\mathrm{Aut}\left(E[p]\right)\mathrm{GL}_2(๐ฝ_p)`$, condition (1) is again seen to characterize the rationality of the point $`\psi ^1(P)`$.
Consider now a non-CM elliptic curve $`E`$ defined over $`k`$ and an isogeny $`\mu :E{}_{}{}^{\nu }E`$ with kernel $`C`$, and assume $`\varrho _E=\varrho `$. This equality amounts to the existence of a basis $`[T_1,T_2]`$ of $`E[p]`$ for which condition (1) holds for every $`\sigma `$ in $`\mathrm{G}_{}`$. We can further suppose that such a basis is sent to either $`\zeta _p`$ or $`\zeta _p^{N^1}`$ by the Weil pairing. In the first case, the point $`P`$ on $`X(N,p)`$ given by the triple $`(E,C,[T_1,T_2]_V)`$ defines through $`\psi `$ a rational point on $`X(N,p)_\varrho `$. In the second case, the triple $`({}_{}{}^{\nu }E,{}_{}{}^{\nu }C,[\mu (T_1),\mu (T_2)]_V)`$ represents a point on $`X(N,p)`$ lying above the same point on $`X^+(N)`$ as $`P`$ and corresponding via $`\psi `$ to a rational point on $`X(N,p)_\varrho `$. Indeed, if we choose the basis $`[\mu (T_1),\mu (T_2)]`$ to fix the isomorphism $`\mathrm{Aut}({}_{}{}^{\nu }E[p])\mathrm{GL}_2(๐ฝ_p)`$, we get the equality $`\varrho _{{}_{}{}^{\nu }E}(\sigma )=\varrho _E(\sigma )`$ for all $`\sigma `$ in $`\mathrm{G}_{}`$.
Lastly, let us consider two different rational points on $`X(N,p)_\varrho `$ corresponding via $`\psi `$ to non-CM points $`P`$ and $`Q`$ on $`X(N,p)`$ with the same image on $`X^+(N)`$. Let the triple $`(E,C,[T_1,T_2]_V)`$ represent the point $`P`$ and fix an isogeny $`\mu :E{}_{}{}^{\nu }E`$ with kernel $`C`$. We must then distinguish two cases for the point $`Q`$ :
* It lies over the pair $`(E,C)`$ on $`X_0(N)`$ if and only if there is a non-trivial element $`\gamma `$ in $`\mathrm{PSL}_2(๐ฝ_p)`$, representing a basis change in $`E[p]`$, such that
$$\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\varrho (\sigma )\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)=\gamma \left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\varrho (\sigma )\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\gamma ^1$$
for all $`\sigma `$ in $`\mathrm{G}_{}`$. This amounts to the existence of a non-trivial element in $`\mathrm{PSL}_2(๐ฝ_p)`$ commuting with all the elements in the image of $`\varrho `$.
* Otherwise, a triple representing $`Q`$ is given by the elliptic curve $`{}_{}{}^{\nu }E`$, the subgroup $`{}_{}{}^{\nu }C`$ and the $`H_V`$โorbit of a basis obtained from $`[\mu (T_1),\mu (T_2)]V`$ by a basis change preserving the Weil pairing. Thus, this case amounts to the existence of an element $`\gamma `$ in $`\mathrm{PSL}_2(๐ฝ_p)`$ such that
$$V\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\varrho (\sigma )\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)V=\gamma \left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\varrho (\sigma )\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\gamma ^1$$
for all $`\sigma `$ in $`\mathrm{G}_{}`$. This is in turn equivalent to the existence of an element in $`\mathrm{PGL}_2(๐ฝ_p)\mathrm{PSL}_2(๐ฝ_p)`$ commuting with all the elements in the image of $`\varrho `$. $`\mathrm{}`$
By the same reasoning as in Section 5, the set of points in Theorem 6.4 is always finite whenever $`N>131`$. A stronger result is obtained from Corollary 3.5.
###### Corollary 6.5
For $`N`$ non-square mod $`p`$, the number of isomorphism classes of $``$-curves of degree $`N`$ realizing $`\varrho `$ is finite, unless $`N=2`$ and $`p=3`$.
## Acknowledgements
I wish to express my gratitude to Joan-Carles Lario for his help and encouragement through this work. I am most indebted to Renรฉ Schoof and the *Dipartimento di Matematica* of the *Universitร di Roma โTor Vergataโ* as well as to the *Dรฉpartement de Mathรฉmatiques de Besanรงon*, where this research was carried out with financial support from the RTN European Network *Galois Theory and Explicit Methods in Arithmetic*. It is also a pleasure to thank Gabriel Cardona for some helpful comments on an earlier version of the paper.
| Julio Fernรกndez |
| --- |
| Departament de Matemร tica Aplicada 4 |
| Universitat Politรจcnica de Catalunya |
| EPSEVG, av. Vรญctor Balaguer |
| E-08800 Vilanova i la Geltrรบ (Barcelona) |
| julio@mat.upc.edu |
|
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# Optical properties of the Q1D multiband models โ the transverse equation of motion approach
## I Introduction
Current investigations of the strongly correlated electron systems often deal with collective contributions to the electrical conductivity (and to the other transport coefficients). This includes in particular the conductivity in the charge-density-wave (CDW) or spin-density-wave (SDW) ordered states. In most of those studies the electrical conductivity is determined, theoretically and experimentally, relative either to the conductivity in the metallic state or with respect to the conductivity in the semiconducting state of the pinned CDW/SDW. Obviously, the prerequisite to such procedures is the accurate evaluation of the conductivity in the metallic state or in the state of the pinned CDW/SDW, and this simple, yet not completely solved, problem is in the focus of our interest here.
Actually, using the microscopic transverse response theory, Lee, Rice and Anderson have found that the single-particle optical conductivity in the ordered CDW state with the negligible small number of scattering centers is given in terms of the semiconducting current-current correlation function which describes excitations across the gap LRA (1, 2, 3). However, it is also shown that for the typical value of the CDW gap and of the zero-frequency damping energy (arising from the impurity scattering processes) their result matches up neither the result of the longitudinal response theory KupcicPB1 (4) nor the experimental observation Degiorgi (5, 6). Yet, the longitudinal and transverse response have to coincide for fast enough (quasi)homogeneous longitudinal fields, as implicit for example in the Maxwell equations of the medium, which employ only one dielectric function.
In his textbook Mahan (7), Mahan has further shown that an alternative, so-called force-force correlation function method gives a good description of the (high-frequency) optical processes in both metallic and semiconducting systems, including various excitations within the conduction band and across the (pseudo)gaps, but also fails to reach correctly the $`\omega 0`$ limit, requiring a specific $`\omega =0`$ field-theory approach Mahan (7, 8, 9).
Also important is the observation that most of the transport-coefficient analyses are based on the Boltzmann equations applied to the nearly-free electron models Mahan (7, 10), completely neglecting band periodicity in the reciprocal space.
Most of these issues can be settled down using the longitudinal response theory with a particular care devoted to the continuity equation Pines (11, 12). In the present article, it will be shown that this can be done alternatively (but in a somewhat less strict way) using the gauge-invariant form of the transverse approach. For this purpose, we consider a quasi-one-dimensional (Q1D) two-band model with the impurity scattering taken into consideration.
The two bands are taken to result from the site-energy dimerization in the highly conducting direction, and, consequently, the Bloch functions and all relevant vertex functions can be determined analytically (Sec. III). Using the equation of motion approach (which is found to be the generalization of the force-force correlation function approach), the intra- and interband optical conductivity are calculated (Sec. IV). In the intraband channel of the transverse correlation functions, the most singular processes in powers of $`1/\omega ^n`$ are collected, resulting in the optical conductivity which matches up the DC conductivity obtained by the Boltzmann equations Ziman (10) or by the Landau response theory Pines (11). The resulting interband conductivity in the CDW ordered state is found to be consistent with the experimental observation. Theorywise, the semiconducting current-current correlation function LRA (1, 2, 3) is replaced by a slightly modified function containing an additional factor which comes from the gauge invariant treatment of the diamagnetic current contributions. Finally (Sec. V), the optical and DC conductivity are determined for a few typical Q1D cases.
## II Transverse multiband response theory
The optical conductivity tensor $`\sigma _{\alpha \alpha }(\omega )`$ is a measure of the absorption rate for the (transverse) electromagnetic waves traveling across the crystal, and the measured spectra, together with the DC conductivity data and other transport coefficients, are an extremely valuable source of the information about the electronic subsystem. Although some aspects of the microscopic response theory can be found in the textbooks Wooten (2, 3, 7, 8, 9, 11, 13), there is no systematic microscopic solution to the multiband optical conductivity problem. Actually, it is easy to determine macroscopic symmetry features of $`\sigma _{\alpha \alpha }(\omega )`$, even in a general case. In this respect, we shall combine the macroscopic symmetry features with the microscopic description of the electron-photon coupling functions (determined for a simple, exactly solvable Q1D electronic model) to develop a consistent microscopic multiband response theory.
### II.1 Optical conductivity tensor
The optical conductivity analysis starts with the Hamiltonian Pines (11)
$`H`$ $`=`$ $`H_0^{\mathrm{field}}+H_0^{\mathrm{el}}+H_1^{\mathrm{ext}}+H_2^{\mathrm{ext}},`$ (1)
which comprises the bare photon term $`H_0^{\mathrm{field}}`$, the bare electronic Hamiltonian $`H_0^{\mathrm{el}}`$ and the first-order and the second-order electron-photon coupling term, $`H_1^{\mathrm{ext}}`$ and $`H_2^{\mathrm{ext}}`$. The bare photon contribution is
$`H_0^{\mathrm{field}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{๐ช\alpha }{}}\left[P_{๐ช\alpha }^{}P_{๐ช\alpha }+\omega _{๐ช0}^2Q_{๐ช\alpha }^{}Q_{๐ช\alpha }\right].`$ (2)
Here $`๐ช`$ and $`\alpha `$ are the wave vector and the polarization of the photon field $`Q_{๐ช\alpha }`$, $`P_{๐ช\alpha }`$ is the field conjugate to $`Q_{๐ช\alpha }`$, and $`\omega _{๐ช0}=cq`$ is the bare photon dispersion. The structure of a typical Q1D tight-binding electronic Hamiltonian $`H^{\mathrm{el}}=H_0^{\mathrm{el}}+H_1^{\mathrm{ext}}+H_2^{\mathrm{ext}}`$ is determined below. However, notice that general symmetry properties of $`\sigma _{\alpha \alpha }(\omega )`$ discussed here do not depend on details of $`H^{\mathrm{el}}`$.
To obtain $`\sigma _{\alpha \alpha }(\omega )`$, the retarded photon Green function is required footnote1 (14)
$`Q_{๐ช\alpha };Q_{๐ช\alpha }_t`$ $`=`$ $`\mathrm{i}\mathrm{\Theta }(t)[Q_{๐ช\alpha }(t),Q_๐ช(0)]`$
$`=`$ $`\mathrm{e}^{\eta t}{\displaystyle \frac{1}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}d\omega \mathrm{e}^{\mathrm{i}\omega t}Q_{๐ช\alpha };Q_{๐ช\alpha }_\omega ,`$
with $`Q_{๐ช\alpha }(t)=\mathrm{e}^{\mathrm{i}Ht/\mathrm{}}Q_{๐ช\alpha }\mathrm{e}^{\mathrm{i}Ht/\mathrm{}}`$ and $`Q_{๐ช\alpha }(0)=Q_{๐ช\alpha }`$ representing the photon fields in the Heisenberg picture at time $`t`$ and $`t=0`$, respectively. Using the equation of motion formalism, we get
$`\left[\omega \left(\omega +\mathrm{i}\eta \right)\omega _{๐ช\alpha }^2\right]Q_{๐ช\alpha };Q_{๐ช\alpha }_\omega `$ $`=`$ $`\mathrm{},`$ (4)
with the adiabatic term $`\eta 0^+`$. The renormalized photon frequency $`\omega _{๐ช\alpha }`$ is given by Pines (11)
$`\omega _{๐ช\alpha }^2`$ $`=`$ $`\omega _{๐ช0}^2+\mathrm{\Omega }_{\mathrm{dia},\alpha }^2+4\pi \mathrm{\Pi }_{\alpha \alpha }(\omega ),`$ (5)
where $`\mathrm{\Omega }_{\mathrm{dia},\alpha }^2`$ and $`4\pi \mathrm{\Pi }_{\alpha \alpha }(\omega )`$ are, respectively, the diamagnetic and current-current contributions to the photon self-energy, shown in Fig. 1 (for the Q1D model under consideration, the explicit form of $`\mathrm{\Pi }_{\alpha \alpha }(\omega )`$ and $`\mathrm{\Omega }_{\mathrm{dia},\alpha }^2`$ is given in Secs. IV B, C and Ref. KupcicPB2 (15), respectively). Combining the Maxwell equations with Eq. (4), it can be shown that the transverse dielectric function $`\epsilon _\alpha (\omega )`$ satisfies the relation
$`\left[\omega \left(\omega +\mathrm{i}\eta \right)\epsilon _\alpha (\omega )\omega _{๐ช0}^2\right]Q_{๐ช\alpha };Q_{๐ช\alpha }_\omega `$ $`=`$ $`\mathrm{},`$ (6)
with
$`\epsilon _\alpha (\omega )`$ $`=`$ $`1+{\displaystyle \frac{4\pi \mathrm{i}}{\omega }}\sigma _{\alpha \alpha }(\omega ).`$ (7)
The optical conductivity defined by Eq. (7) is
$`\sigma _{\alpha \alpha }(\omega )`$ $`=`$ $`{\displaystyle \frac{\mathrm{i}}{4\pi \left(\omega +\mathrm{i}\eta \right)}}\left[\mathrm{\Omega }_{\mathrm{dia},\alpha }^2+4\pi \mathrm{\Pi }_{\alpha \alpha }(\omega )\right],`$ (8)
irrespective of whether the electromagnetic field is treated as a classical or a quantum field. This is a quite general and in many respects very useful result. In the general case, with several valence bands or with several scattering channels in a single band, $`\mathrm{\Omega }_{\mathrm{dia},\alpha }^2`$ includes one or more diamagnetic contributions (depending on the number of bands intersecting the Fermi level), and $`\mathrm{\Pi }_{\alpha \alpha }(\omega )`$ represents all intra- and interband current-current correlation functions.
The following general properties of the expression (8) are important for interpreting the measured spectra. First, there are at least two distinct structures in the optical conductivity spectrum $`\mathrm{Re}\{\sigma _{\alpha \alpha }(\omega )\}`$, the first one is a delta function at $`\omega =0`$, related to the diamagnetic current, and the second one represents various contributions, including the exciton contributions if the short-range dipole-dipole interactions are present Wooten (2, 16, 17, 18). This is easily seen from
$`\mathrm{Re}\{\sigma _{\alpha \alpha }(\omega )\}`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left[\mathrm{\Omega }_{\mathrm{dia},\alpha }^2+4\pi \mathrm{Re}\{\mathrm{\Pi }_{\alpha \alpha }(0)\}\right]\delta (\omega )`$ (9)
$`{\displaystyle \frac{1}{\omega }}\mathrm{Im}\{\mathrm{\Pi }_{\alpha \alpha }(\omega )\}.`$
Obviously, in the normal metallic or insulating state the delta-function term vanishes Kittel (13). Therefore, any consistent treatment of the electron-photon coupling functions has to fulfill the relation
$`\mathrm{\Omega }_{\mathrm{dia},\alpha }^2+4\pi \mathrm{Re}\{\mathrm{\Pi }_{\alpha \alpha }(0)\}=0.`$ (10)
The total optical conductivity in the normal metallic or insulating state can be written then in the form
$`\sigma _{\alpha \alpha }(\omega )`$ $`=`$ $`{\displaystyle \frac{\mathrm{i}}{\omega }}\left[\mathrm{\Pi }_{\alpha \alpha }(\omega )\mathrm{Re}\{\mathrm{\Pi }_{\alpha \alpha }(0)\}\right].`$ (11)
Noteworthy, in the superconducting state the sum
$`{\displaystyle \frac{1}{2}}\left[\mathrm{\Omega }_{\mathrm{dia},\alpha }^2+4\pi \mathrm{Re}\{\mathrm{\Pi }_{\alpha \alpha }(0)\}\right]`$ (12)
measures the weight of the missing area in optical conductivity spectra Schrieffer (19, 20, 21), while the single-particle contributions are still given by Eq. (11). Finally, notice that in absence of local dipolar excitations (the case of the site-energy dimerization discussed in this article) the total current-current correlation function is the sum of only two contributions describing, respectively, the creation of the โfreeโ intra- and interband electron-hole pairs, while the processes associated with excitons (i.e. the quasiparticles representing the bound electron-hole pairs) are not present.
The optical conductivity determined by Eq. (11), with the current-current correlation function calculated by the equation of motion approach, is in the focus of the present analysis. The causality requirement Wooten (2, 13) (i.e. the KramersโKronig relations), the effective mass theorem KupcicPB1 (4, 22, 23) and the gauge-invariance requirement KupcicPB1 (4, 11, 15) will be used to test the obtained results. A particular care will be devoted to the non-physical singularity at $`\omega =0`$ related to the prefactor of $`\omega ^1`$ in Eq. (11), and to the construction of the optical conductivity model with the correct behaviour in the $`\omega 0`$ limit, giving rise to a unified description of the optical and transport phenomena.
## III Electronic Hamiltonian with two qualitatively different scattering channels
Although our model is Q1D, the present response theory is quite general and the electronic Hamiltonian could represent an arbitrary multiband model. We assume that in addition to the bare electronic Hamiltonian (denoted below by $`H_0`$) there is a static single-particle potential ($`H_0^{}`$) characterized by a commensurate wave vector, which describes, for example, the scattering of electrons on the site-energy dimerization potential. The other single-electron scattering processes (on impurities, phonons, etc.) are represented by $`H_1^{}`$. For most of the questions discussed here the main effects of the two-electron interactions are taken satisfactorily into account through the effective mean fields included in $`H_0^{}`$ or $`H_1^{}`$, footnote2 (24) resulting finally in $`H_0^{\mathrm{el}}=H_0+H_0^{}+H_1^{}`$ in Eq. (1). $`H^{\mathrm{ext}}=H_1^{\mathrm{ext}}+H_2^{\mathrm{ext}}`$ couples the valence electrons to transverse electromagnetic fields. The resulting total electronic Hamiltonian is
$`H^{\mathrm{el}}`$ $`=`$ $`H_0^{\mathrm{el}}+H^{\mathrm{ext}}.`$ (13)
We start by diagonalizing the Hamiltonian $`H_0+H_0^{}`$. As mentioned above, we consider the simplest Q1D model, where $`H_0^{}`$ represents the site-energy dimerization in the highly conducting direction and in $`H_1^{}`$ only the impurity scattering is taken into account. Next, we determine the related electron-photon coupling functions. At the end of this section, the multiband current-current correlation function is introduced.
### III.1 Bare Hamiltonian
The single-particle properties of the Q1D site-energy-dimerization model come from the exact diagonalization of the Hamiltonian KupcicPB2 (15)
$`H_0+H_0^{}`$ $`=`$ $`{\displaystyle \underset{๐ค\sigma }{}}[\epsilon _c(๐ค)c_{๐ค\sigma }^{}c_{๐ค\sigma }+\epsilon _{\underset{ยฏ}{c}}(๐ค)\underset{ยฏ}{c}_{๐ค\sigma }^{}\underset{ยฏ}{c}_{๐ค\sigma }`$ (14)
$`+\mathrm{\Delta }(\underset{ยฏ}{c}_{๐ค\sigma }^{}c_{๐ค\sigma }+c_{๐ค\sigma }^{}\underset{ยฏ}{c}_{๐ค\sigma })].`$
The bare electron dispersions of two subbands, artificially dimerized along the highly conducting direction $`a`$, are
$`\epsilon _{\underset{ยฏ}{c},c}(๐ค)`$ $`=`$ $`\pm 2t_a\mathrm{cos}๐ค๐2t_b\mathrm{cos}๐ค๐,`$ (15)
with $`\epsilon _{\underset{ยฏ}{c}}(๐ค)\epsilon _c(๐ค\pm \pi /a\widehat{x})`$, and with the wave vector k restricted to the new (reduced) Brillouin zone, $`0.5\pi /ak_x0.5\pi /a`$, $`\pi /bk_y\pi /b`$. $`t_a`$ and $`t_b`$ ($`t_at_b>0`$) are the bond energies in the direction $`a`$ and the perpendicular direction $`b`$, respectively, and $`\mathrm{\Delta }`$ is the magnitude of the dimerization potential in the direction $`a`$. Note that such potential corresponds to imperfect nesting in Eq. (15) in contrast to dimerization in all directions; the former nesting is chosen because it is more interesting in context of the conductivity studies.
The transformations between the unperturbed states, the band index $`l=c,\underset{ยฏ}{c}`$, and the Bloch states, the band index $`L=C,\underset{ยฏ}{C}`$, of the form
$`l_{๐ค\sigma }^{}`$ $`=`$ $`{\displaystyle \underset{L}{}}U_๐ค(l,L)L_{๐ค\sigma }^{},`$ (16)
lead to
$`H_0`$ $`=`$ $`{\displaystyle \underset{L๐ค\sigma }{}}E_L(๐ค)L_{๐ค\sigma }^{}L_{๐ค\sigma },`$ (17)
with the dispersions
$`E_{\underset{ยฏ}{C},C}(๐ค)`$ $`=`$ $`{\displaystyle \frac{1}{2}}[\epsilon _{\underset{ยฏ}{c}}(๐ค)+\epsilon _c(๐ค)]\pm \sqrt{{\displaystyle \frac{1}{4}}\epsilon _{\underset{ยฏ}{c}c}^2(๐ค)+\mathrm{\Delta }^2},`$
$`\epsilon _{\underset{ยฏ}{c}c}(๐ค)`$ $`=`$ $`\epsilon _{\underset{ยฏ}{c}}(๐ค)\epsilon _c(๐ค).`$ (18)
The transformation-matrix elements are given by
$`\left(\begin{array}{cc}U_๐ค(c,C)\hfill & U_๐ค(c,\underset{ยฏ}{C})\hfill \\ U_๐ค(\underset{ยฏ}{c},C)\hfill & U_๐ค(\underset{ยฏ}{c},\underset{ยฏ}{C})\hfill \end{array}\right)`$ $`=`$ $`\left(\begin{array}{cc}u(๐ค)& v(๐ค)\\ v(๐ค)& u(๐ค)\end{array}\right),`$ (23)
where $`u(๐ค)=\mathrm{cos}{\displaystyle \frac{\phi (๐ค)}{2}}`$, $`v(๐ค)=\mathrm{sin}{\displaystyle \frac{\phi (๐ค)}{2}}`$, with the auxiliary phase $`\phi (๐ค)`$ defined in the usual way
$`\mathrm{tan}\phi (๐ค)`$ $`=`$ $`{\displaystyle \frac{2\mathrm{\Delta }}{\epsilon _{\underset{ยฏ}{c}c}(๐ค)}}.`$ (24)
The bands are shown in Fig. 2. Hereafter, the lower (conduction) band is assumed to be partially filled and the upper (valence) band is empty. For the lower band completely filled and $`\mathrm{\Delta }t_b`$, the band structure corresponds to the commensurate CDW system, otherwise we have the metallic behaviour.
### III.2 Electron-photon coupling Hamiltonian
Using the generalized minimal substitution method for the tight-binding electrons KupcicPB1 (4, 15), we obtain that the conduction electrons described by the Hamiltonian (14) are coupled to the external electromagnetic fields through
$`H^{\mathrm{ext}}`$ $`=`$ $`{\displaystyle \underset{l}{}}{\displaystyle \underset{๐ค\sigma }{}}{\displaystyle \underset{๐ช\alpha }{}}\delta H_0^l(๐ค,๐ช)l_{๐ค+๐ช\sigma }^{}l_{๐ค\sigma },`$ (25)
where, to the second order in the vector potential $`A_\alpha (๐ช)`$,
$`\delta H_0^l(๐ค,๐ช)`$ $``$ $`{\displaystyle \frac{\epsilon _l(๐ค)}{k_\alpha }}{\displaystyle \frac{e}{\mathrm{}c}}A_\alpha (๐ช)+{\displaystyle \frac{1}{2}}{\displaystyle \frac{^2\epsilon _l(๐ค)}{k_\alpha ^2}}\left({\displaystyle \frac{e}{\mathrm{}c}}\right)^2A_\alpha ^2(๐ช).`$
The photon annihilation operator $`A_{๐ช\alpha }`$ enters in Eq. (LABEL:eq22) through
$`A_\alpha (๐ช)=\sqrt{{\displaystyle \frac{4\pi c^2}{V}}}Q_{๐ช\alpha },`$
$`A_\alpha ^2(๐ช)={\displaystyle \underset{๐ช^{}}{}}A_\alpha (๐ช๐ช^{})A_\alpha (๐ช^{}),`$
where
$`Q_{๐ช\alpha }`$ $`=`$ $`\sqrt{{\displaystyle \frac{\mathrm{}}{2\omega _{๐ช\alpha }}}}\left[A_{๐ช\alpha }+A_{๐ช\alpha }^{}\right]`$ (27)
is the electromagnetic displacement field of Eq. (2) Pines (11). Finally, in the Bloch representation the coupling Hamiltonian becomes
$`H^{\mathrm{ext}}`$ $`=`$ $`H_1^{\mathrm{ext}}+H_2^{\mathrm{ext}}={\displaystyle \frac{1}{c}}{\displaystyle \underset{๐ช\alpha }{}}A_\alpha (๐ช)\widehat{J}_\alpha (๐ช)`$ (28)
$`+{\displaystyle \frac{e^2}{2mc^2}}{\displaystyle \underset{๐ช\alpha }{}}A_\alpha ^2(๐ช)\widehat{\gamma }_{\alpha \alpha }(๐ช;2),`$
with
$`\widehat{J}_\alpha (๐ช)`$ $`=`$ $`{\displaystyle \underset{LL^{}}{}}{\displaystyle \underset{๐ค\sigma }{}}J_\alpha ^{LL^{}}(๐ค)L_{๐ค\sigma }^{}L_{๐ค+๐ช\sigma }^{},`$ (29)
$`\widehat{\gamma }_{\alpha \alpha }(๐ช;2)`$ $`=`$ $`{\displaystyle \underset{LL^{}}{}}{\displaystyle \underset{๐ค\sigma }{}}\gamma _{\alpha \alpha }^{LL^{}}(๐ค;2)L_{๐ค\sigma }^{}L_{๐ค+๐ช\sigma }^{},(25^{})`$
representing, respectively, the current density and Raman Genkin (22, 15) density operators. The structure of the related vertex functions $`J_\alpha ^{LL^{}}(๐ค)`$ and $`\gamma _{\alpha \alpha }^{LL^{}}(๐ค;2)`$ is given in Appendix A.
### III.3 Single-particle scattering processes
We consider here only the intraband impurity scattering processes in the perturbation
$`H_1^{}`$ $`=`$ $`{\displaystyle \underset{L}{}}{\displaystyle \underset{\mathrm{๐ค๐ค}^{}\sigma }{}}V^{LL}(๐ค,๐ค^{})L_{๐ค\sigma }^{}L_{๐ค^{}\sigma },`$ (30)
for which one usually assumes $`V^{LL}(๐ค,๐ค^{})=V^{LL}(๐ค๐ค^{})`$ Abrikosov (8, 9). This form of $`H_1^{}`$ is consistent with the regime in which $`\epsilon _{\underset{ยฏ}{c}c}(๐ค_\mathrm{F})\mathrm{}/\tau `$ ($`\tau `$ is the relaxation time defined below). The generalization is straightforward.
### III.4 Current-current correlation function
In the equation of motion approach Pines (11) used in the next section, the starting point is the current-current correlation function of Eqs. (8)โ(12) shown in Fig. 3. It is defined by LRA (1, 2, 7, 15)
$`\mathrm{\Pi }_{\alpha \alpha }(๐ช,t)`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{}V}}\widehat{J}_\alpha (๐ช);\widehat{J}_\alpha (๐ช)_t`$ (31)
$``$ $`{\displaystyle \frac{\mathrm{i}}{\mathrm{}V}}\mathrm{\Theta }(t)[\widehat{J}_\alpha (๐ช,t),\widehat{J}_\alpha (๐ช,0)].`$
Here the current operator $`\widehat{J}_\alpha (๐ช)`$ includes all intra- and interband current density fluctuations, as seen from Eq. (29). According to Fig. 3, $`\mathrm{\Pi }_{\alpha \alpha }(๐ช,t)`$ comprises two intraband and two interband contributions, and the problem is reduced, as will be seen immediately below, to the self-consistent calculation of intra- and interband electron-hole propagator in presence of the perturbation $`H_1^{}`$. In Eq. (31), as well as in Eqs. (32), (34), (37) and (38), the usual notation for the retarded correlation functions is used: $`\widehat{A},\widehat{B}_t=\mathrm{i}\mathrm{\Theta }(t)[\widehat{A}(t),\widehat{B}(0)]`$, with $`\widehat{A}(t)`$ being the operator $`\widehat{A}`$ in the Heisenberg picture.
## IV Equation of motion approach
### IV.1 Generalized correlation functions
The retarded electron-hole propagator
$`๐_1^{LL^{}}(๐ค,๐ค_+,๐ค_+^{},๐ค^{},t)=L_{๐ค\sigma }^{}L_{๐ค+๐ช\sigma }^{};L_{}^{}{}_{๐ค^{}+๐ช\sigma }{}^{}L_{๐ค^{}\sigma }_t`$
$`\mathrm{i}\mathrm{\Theta }(t)[\left(L_{๐ค\sigma }^{}L_{๐ค+๐ช\sigma }^{}\right)_t,\left(L_{}^{}{}_{๐ค^{}+๐ช\sigma }{}^{}L_{๐ค^{}\sigma }\right)_0]`$
$`=\mathrm{e}^{\eta t}{\displaystyle \frac{1}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}d\omega \mathrm{e}^{\mathrm{i}\omega t}๐_1^{LL^{}}(๐ค,๐ค_+^{},๐ค_+,๐ค^{},\omega )`$ (32)
($`๐ค_+`$ is the abbreviation for $`๐ค+๐ช`$) is the central quantity to all long-wavelength correlation functions, as can be seen from Fig. 4, or from the expression
$`\chi _{f,g}(๐ช,t)={\displaystyle \frac{1}{\mathrm{}V}}{\displaystyle \underset{LL^{}}{}}{\displaystyle \underset{\mathrm{๐ค๐ค}^{}\sigma }{}}f^{LL^{}}(๐ค,๐ค+๐ช)g^{L^{}L}(๐ค^{}+๐ช,๐ค^{})`$
$`\times ๐_1^{LL^{}}(๐ค,๐ค_+,๐ค_+^{},๐ค^{},t),`$ (33)
which represents a generalized long-wavelength correlation function, with $`f^{LL^{}}(๐ค,๐ค+๐ช)`$ and $`g^{L^{}L}(๐ค^{}+๐ช,๐ค^{})`$ being the charge, current, Raman or some other vertex functions. In Eq. (32) the abbreviation $`(\widehat{A}\widehat{B})_t=\widehat{A}(t)\widehat{B}(t)`$ is used.
The symmetry of the $`f^{LL^{}}(๐ค,๐ค+๐ช)`$ and $`g^{L^{}L}(๐ค^{}+๐ช,๐ค^{})`$ vertices, together with the nature of the singularity of the leading term in the perturbation series, determines the summation rule for the related Feynman diagrams. The simplest way to collect the most singular diagrams is to consider the equations of motion connecting $`๐_1^{LL^{}}(๐ค,๐ค_+,๐ค_+^{},๐ค^{},t)`$ with the correlation function $`๐_2^{LL^{}}(๐ค,๐ค_+,๐ค_+^{},๐ค^{},t)`$ defined as
$`๐_2^{LL^{}}(๐ค,๐ค_+,๐ค_+^{},๐ค^{},t)`$ (34)
$`=[L_{๐ค\sigma }^{}L_{๐ค+๐ช\sigma }^{},H^{}];L_{}^{}{}_{๐ค^{}+๐ช\sigma }{}^{}L_{๐ค^{}\sigma }_t.`$
The direct calculation gives the exact relation
$`\mathrm{}\left[๐_0^{LL^{}}(๐ค,๐ค_+,\omega )\right]^1๐_1^{LL^{}}(๐ค,๐ค_+,๐ค_+^{},๐ค^{},\omega )`$ (35)
$`=\mathrm{}\delta _{๐ค,๐ค^{}}\left[f_L(๐ค)f_L^{}(๐ค_+)\right]+๐_2^{LL^{}}(๐ค,๐ค_+,๐ค_+^{},๐ค^{},\omega ).`$
Here
$`\mathrm{}\left[๐_0^{LL^{}}(๐ค,๐ค_+,\omega )\right]^1=\mathrm{}(\omega +\mathrm{i}\eta )+E_L(๐ค)E_L^{}(๐ค_+)`$
(36)
is a useful abbreviation. $`๐_2^{LL^{}}(๐ค,๐ค_+,๐ค_+^{},๐ค^{},\omega )`$ is the Fourier transform of $`๐_2^{LL^{}}(๐ค,๐ค_+,๐ค_+^{},๐ค^{},t)`$, and $`f_L(๐ค)f(E_L(๐ค))=L_{๐ค\sigma }^{}L_{๐ค\sigma }=\left[1+\mathrm{exp}\{\beta [E_L(๐ค)\mu ]\}\right]^1`$ is the FermiโDirac function, with $`\beta =1/(k_\mathrm{B}T)`$.
The way to evaluate $`๐_2^{LL^{}}(๐ค,๐ค_+,๐ค_+^{},๐ค^{},t)`$ depends on the choice of representation of this function. There are two alternative ways,
$`๐_2^{LL^{}}(๐ค,๐ค_+,๐ค_+^{},๐ค^{},t)`$ (37)
$`=\mathrm{i}\mathrm{\Theta }(t)[[L_{๐ค\sigma }^{}L_{๐ค+๐ช\sigma }^{},H^{}]_t,\left(L_{}^{}{}_{๐ค^{}+๐ช\sigma }{}^{}L_{๐ค^{}\sigma }\right)_0]`$
or
$`๐_2^{LL^{}}(๐ค,๐ค_+,๐ค_+^{},๐ค^{},t)`$ (38)
$`=\mathrm{i}\mathrm{\Theta }(t)[[L_{๐ค\sigma }^{}L_{๐ค+๐ช\sigma }^{},H^{}]_0,\left(L_{}^{}{}_{๐ค^{}+๐ช\sigma }{}^{}L_{๐ค^{}\sigma }\right)_t],`$
leading to two different self-consistent schemes (one described below and another encountered in the longitudinal response theory Pines (11)), both giving, as will be argued below, the same result. Again, $`[\widehat{A},\widehat{B}]_t`$ is the abbreviation for $`[\widehat{A}(t),\widehat{B}(t)]`$.
It is important to realize at the outset that the first term on the right-hand side of Eq. (35) is significant for all interband correlation functions $`\chi _{f,g}(๐ช,t)`$ in Eq. (33), irrespective of the vertex symmetries. The probability for the direct creation of the interband electron-hole pair is proportional here to $`f_L(๐ค)f_L^{}(๐ค_+)f_L(๐ค)f_L^{}(๐ค)`$, so that all occupied states in the conduction band(s) are equally important. For vertices $`f^{LL^{}}(๐ค,๐ค+๐ช)`$ and $`g^{L^{}L}(๐ค^{}+๐ช,๐ค^{})`$ taken to represent the intraband current vertices, the correlation function $`\chi _{f,g}(๐ช,t)`$ in Eq. (33) becomes the intraband current-current correlation function $`\mathrm{\Pi }_{\alpha \alpha }^{\mathrm{intra}}(๐ช,t)`$. According to the discussion of Sec. II A, the related intraband optical conductivity is ruled by the prefactor $`\omega ^1`$ in Eq. (11). The direct processes in Eq. (35) (see the first term on the right-hand side of Fig. 3(c)) related to $`f_C(๐ค)f_C(๐ค_+)0`$ are insignificant, and $`\mathrm{\Pi }_{\alpha \alpha }^{\mathrm{intra}}(๐ช,t)`$ can be adequately described by the second, indirect term in this equation (the second term in Fig. 3(c); see also Fig. 5), with $`๐_2^{LL^{}}(๐ค,๐ค_+,๐ค_+^{},๐ค^{},t)`$ given by Eq. (38).
Before turning to the evaluation of $`\mathrm{\Pi }_{\alpha \alpha }^{\mathrm{intra}}(๐ช,t)`$ it is interesting to contrast this conclusion with its analog in the longitudinal response theory. In the longitudinal approach, the long-range Coulomb forces are activated, and the first term in Eq. (35) dominates the intraband charge-charge correlation function, even in the dynamic limit, since the small probability for the intraband electron-hole pair creation, proportional to $`f_C(๐ค)f_C(๐ค_+)`$, cancels out the $`q^2`$ singularity of the long-range forces, which is the well-known RPA result. The indirect scattering processes, on the other hand, are proportional to the effective intraband charge vertex, analogous to the effective intraband current vertex introduced below, Eq. (45). Since the bare long-wavelength intraband charge vertex satisfies $`q(๐ค+๐ช,๐ค)q(๐ค,๐ค)=e`$, where $`e`$ is the bare electron charge, the effective intraband charge vertex for the indirect scattering processes vanishes, because $`q(๐ค,๐ค)q(๐ค^{},๐ค^{})=0`$. Since the longitudinal and transverse approaches are to be equivalent, this means that the direct longitudinal processes are to be equivalent to the indirect transverse processes in the intraband channel, while the contributions of both the indirect longitudinal and the direct transverse scattering processes are to be negligible. This issue will be further discussed in Sec. IV B 5.
The above transverse approach can now be applied to the current-current correlation function of the two band model, rewritten in the form
$`\mathrm{\Pi }_{\alpha \alpha }(๐ช,t)`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{}V}}{\displaystyle \underset{LL^{}}{}}{\displaystyle \underset{\mathrm{๐ค๐ค}^{}\sigma }{}}J_\alpha ^{LL^{}}(๐ค)J_\alpha ^{L^{}L}(๐ค^{})`$ (39)
$`\times ๐_1^{LL^{}}(๐ค,๐ค_+,๐ค_+^{},๐ค^{},t).`$
Here the vertices $`f^{LL^{}}(๐ค,๐ค+๐ช)`$ and $`g^{L^{}L}(๐ค^{}+๐ช,๐ค^{})`$ in Eq. (33) are replaced by $`J_\alpha ^{LL^{}}(๐ค,๐ค+๐ช)J_\alpha ^{LL^{}}(๐ค)`$ and $`J_\alpha ^{L^{}L}(๐ค^{}+๐ช,๐ค^{})J_\alpha ^{L^{}L}(๐ค^{})`$, respectively. Optical processes relevant to the two-band model (including the indirect interband ones not considered in the present analysis) are illustrated in Fig. 5.
In order to make presentation of the results more transparent, we shall first determine the intraband contributions, and then give the analysis of the interband optical excitations. For the sake of brevity, in the next section the (intra)band index $`C`$ will be omitted ($`E_C(๐ค)E(๐ค)`$, $`J_\alpha ^{CC}(๐ค)J_\alpha (๐ค)`$, โฆ, with $`C_{๐ค\sigma }^{}c_{๐ค\sigma }^{}`$). It should be noticed that the results obtained below for the intraband optical conductivity are quite general, i.e. they cover various physically different regimes. As explained in Sec. IV B 2, when the electron filling of the dimerized band varies between 0 and 1, the electronic system transforms from an electron-like semiconducting, through a metallic, into a hole-like semiconducting regime. A more detailed discussion of this issue is given in Sec. V, in context of the total optical conductivity.
### IV.2 Intraband optical conductivity
According to the aforementioned arguments, the intraband optical processes are described by the equations
$`\mathrm{}๐_0^1(๐ค,๐ค_+,\omega )๐_1(๐ค,๐ค_+,๐ค_+^{},๐ค^{},\omega )`$
$`=๐_2(๐ค,๐ค_+,๐ค_+^{},๐ค^{},\omega ),`$ (40)
$`\mathrm{}๐_0^1(๐ค^{},๐ค_+^{},\omega )๐_2(๐ค,๐ค_+,๐ค_+^{},๐ค^{},\omega )`$
$`=๐_3(๐ค,๐ค_+,๐ค_+^{},๐ค^{},\omega )`$
$`+\mathrm{}[[c_{๐ค\sigma }^{}c_{๐ค+๐ช\sigma },H^{}],c_{๐ค^{}+๐ช\sigma }^{}c_{๐ค^{}\sigma }],`$ (41)
where $`๐_0^1(๐ค^{},๐ค_+^{},\omega )`$ is the intraband term in Eq. (36) and $`๐_3(๐ค,๐ค_+,๐ค_+^{},๐ค^{},\omega )`$ is the Fourier transform of the force-force correlation function Mahan (7, 25, 15)
$`๐_3(๐ค,๐ค_+,๐ค_+^{},๐ค^{},t)`$ (42)
$`=\mathrm{i}\mathrm{\Theta }(t)[[c_{๐ค\sigma }^{}c_{๐ค+๐ช\sigma },H^{}]_0,[c_{๐ค^{}+๐ช\sigma }^{}c_{๐ค^{}\sigma },H^{}]_t].`$
The second term in Eq. (41) is the ground-state average of the four-operator product at $`t=0`$. It is off-diagonal in the Bloch representation, and its value can be obtained by putting $`E(๐ค_+^{})E(๐ค^{})0`$ in the left-hand side of Eq. (41) and then taking the formal limit $`\omega 0`$. The result is
$`\mathrm{}[[c_{๐ค\sigma }^{}c_{๐ค+๐ช\sigma },H^{}],c_{๐ค^{}+๐ช\sigma }^{}c_{๐ค^{}\sigma }]`$
$`๐_3(๐ค,๐ค_+,๐ค_+^{},๐ค^{},\omega =0).`$ (43)
For the impurity scattering, we have the cancellation of this constant term with its counterpart obtained by the electron $``$ hole replacement (see Eq. (51), for example), but for a time-dependent perturbation $`H_1^{}`$, Eq. (43) achieves an interesting structure, as pointed out in Refs. Gotze (25, 26)
Due to the large velocity of light, the energy and wave vector transfers in the external points of the intraband current-current correlation function fulfill $`\mathrm{}\omega E(๐ค_+^{})E(๐ค^{})`$, and the factors $`๐_0(๐ค,๐ค_+,\omega )`$ and $`๐_0(๐ค^{},๐ค_+^{},\omega )`$ in Eqs. (40) and (41), representing the propagator of the virtual electron-hole pairs (related to the process $`13`$ in Fig. 5), can be replaced by $`1/\omega `$. The resulting intraband correlation function becomes
$`\mathrm{\Pi }_{\alpha \alpha }^{\mathrm{intra}}(\omega )={\displaystyle \frac{1}{\mathrm{}V(\mathrm{}\omega )^2}}{\displaystyle \underset{\mathrm{๐ค๐ค}^{\prime \prime }๐ค_1๐ค_1^{\prime \prime }\sigma }{}}\left[J_\alpha (๐ค)J_\alpha (๐ค^{\prime \prime })\right]`$
$`\times [J_\alpha (๐ค_1^{\prime \prime })J_\alpha (๐ค_1)]<V(๐ค๐ค^{\prime \prime })V(๐ค_1^{\prime \prime }๐ค_1)`$
$`\times [๐_1(๐ค,๐ค^{\prime \prime },๐ค_1^{\prime \prime },๐ค_1,\omega )๐_1(๐ค,๐ค^{\prime \prime },๐ค_1^{\prime \prime },๐ค_1,0)]>.`$ (44)
$`\mathrm{}`$ is here and subsequently the average over the impurity sites Mahan (7).
The diagrammatic representation of $`\mathrm{\Pi }_{\alpha \alpha }^{\mathrm{intra}}(๐ช,\omega )`$ is given in Fig. 6(a), with the structure of the leading term, proportional to $`(H_1^{})^2/\omega `$, shown explicitly in Fig. 6(b). The expression
$`j_\alpha (๐ค,๐ค^{\prime \prime })`$ $`=`$ $`{\displaystyle \frac{V(๐ค๐ค^{\prime \prime })}{\mathrm{}\omega }}\left[J_\alpha (๐ค)J_\alpha (๐ค^{\prime \prime })\right]`$ (45)
can be recognized as an effective current vertex for the indirect intraband photon absorption/emission processes, and is a sum of two terms encircled in Fig. 6(b). Notice also that $`๐_1(๐ค,๐ค^{\prime \prime },๐ค_1^{\prime \prime },๐ค_1,\omega )๐_1(๐ค,๐ค^{\prime \prime },๐ค_1^{\prime \prime },๐ค_1,0)`$ is proportional to $`\omega `$, canceling out the factor of $`\omega ^1`$ in one of the effective vertices.
It is important to notice that in the leading term (Fig. 6(b)) the expression (44) is identical to the result of the force-force correlation function method Mahan (7). While the latter method is usually limited to the examination of this leading term, or to the summation of irrelevant higher order diagrams which results in the well-known Hopfield formula Mahan (7), the present approach is focused instead on the exact summation of the most singular contributions in powers of $`(H_1^{})^2/\omega `$ and should be regarded as a generalization of the standard force-force correlation function approach.
#### IV.2.1 Proper electron-hole representation
After determining the structure of the external points (i.e. the effective vertices) in the diagram shown in Fig. 6(a), we have to find the internal structure of the diagram. The latter is represented by the indirect (impurity-assisted) electron-hole propagator $`๐_1(๐ค,๐ค^{\prime \prime },๐ค_1^{\prime \prime },๐ค_1,\omega )`$, characterized by the momentum transfer $`๐ค๐ค^{\prime \prime }`$, rather than by the negligibly small external momentum transfer $`๐ค_+๐ค=๐ช`$ of the direct processes in Eq. (35). Its internal structure is determined here by the self-consistent solution of the exact equation
$`\mathrm{}๐_0^1(๐ค,๐ค^{\prime \prime },\omega )๐_1(๐ค,๐ค^{\prime \prime },๐ค_1^{\prime \prime },๐ค_1,\omega )`$ (46)
$`=\mathrm{}\delta _{๐ค,๐ค_1}\delta _{๐ค^{\prime \prime },๐ค_1^{\prime \prime }}\left[f(๐ค)f(๐ค^{\prime \prime })\right]+๐_2(๐ค,๐ค^{\prime \prime },๐ค_1^{\prime \prime },๐ค_1,\omega ).`$
$`๐_1(๐ค,๐ค^{\prime \prime },๐ค_1^{\prime \prime },๐ค_1,\omega )`$ represents the electron-hole pair created by the indirect photon absorption of Fig. 5, which means that the energy transfer $`\mathrm{}\omega `$ is close to the electron-hole pair energy $`E(๐ค)E(๐ค^{\prime \prime })`$. The wave vectors k and $`๐ค^{\prime \prime }`$ ($`๐ค_1`$ and $`๐ค_1^{\prime \prime }`$, as well) are independent of each other, and therefore the first term in Eq. (46) dominates the low-energy physics, preferring the optical processes between the states $`E(๐ค)\mu `$ and $`E(๐ค^{\prime \prime })\mu `$.
In combining Eq. (46) with the equation of motion for $`๐_2(๐ค,๐ค^{\prime \prime },๐ค_1^{\prime \prime },๐ค_1,\omega )`$ (see Eq. (57) for more detail) it is sufficient to collect only the terms in the summation which are relevant to the self-consistent equation for the kernel $`_\alpha (๐ค_1^{\prime \prime },๐ค_1,\omega )`$ in the current-current correlation function (see Figs. 7(c,d), and the criterion of the validity shown in Figs. 7(a,b))
$`\mathrm{\Pi }_{\alpha \alpha }^{\mathrm{intra}}(\omega )={\displaystyle \frac{1}{\mathrm{}V(\mathrm{}\omega )^2}}{\displaystyle \underset{๐ค_1๐ค_1^{\prime \prime }\sigma }{}}<\left[J_\alpha (๐ค_1^{\prime \prime })J_\alpha (๐ค_1)\right]`$
$`\times [_\alpha (๐ค_1^{\prime \prime },๐ค_1,\omega )_\alpha (๐ค_1^{\prime \prime },๐ค_1,0)]>,`$ (47)
where
$`_\alpha (๐ค_1^{\prime \prime },๐ค_1,\omega )={\displaystyle \underset{\mathrm{๐ค๐ค}^{\prime \prime }}{}}V(๐ค๐ค^{\prime \prime })V(๐ค_1^{\prime \prime }๐ค_1)`$
$`\times \left[J_\alpha (๐ค)J_\alpha (๐ค^{\prime \prime })\right]๐_1(๐ค,๐ค^{\prime \prime },๐ค_1^{\prime \prime },๐ค_1,\omega ).`$ (48)
The idea of the present approach is the self-consistent treatment of Eqs. (46) and (IV.2.1), which represent two equations connecting the electron-hole propagator $`๐_1(๐ค,๐ค^{\prime \prime },๐ค_1^{\prime \prime },๐ค_1,\omega )`$ with the kernel $`_\alpha (๐ค_1^{\prime \prime },๐ค_1,\omega )`$. The solution is based on the expansion
$`๐_1(๐ค,๐ค^{\prime \prime },๐ค_1^{\prime \prime },๐ค_1,\omega )={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}๐_1^{(2n)}(๐ค,๐ค^{\prime \prime },๐ค_1^{\prime \prime },๐ค_1,\omega ),`$ (49)
where $`๐_1^{(2n)}(๐ค,๐ค^{\prime \prime },๐ค_1^{\prime \prime },๐ค_1,\omega )`$ includes only the self-consistent terms proportional to $`(H_1^{})^{2n}/\omega ^n`$.
As easily seen in the longitudinal analysis, the self-consistent expression for $`๐_1(๐ค,๐ค^{\prime \prime },๐ค_1^{\prime \prime },๐ค_1,\omega )`$, with $`๐ค^{\prime \prime }๐ค+๐ช`$, plays the leading role. Accordingly, to compare both response theories, one needs the self-consistent scheme for $`_\alpha (๐ค_1^{\prime \prime },๐ค_1,\omega )`$ and the recurrence relations for $`๐_1^{(2n)}(๐ค,๐ค^{\prime \prime },๐ค_1^{\prime \prime },๐ค_1,\omega )`$ (see Sec. IV B 3).
For pedagogical reasons, it is convenient first to consider the zeroth order contribution to (47) and define the effective number of conduction electrons and the related electron-hole damping energy.
#### IV.2.2 High-frequency limit
The direct calculation gives the first term in the expansion (49)
$`๐_1^{(0)}(๐ค,๐ค^{\prime \prime },๐ค_1^{\prime \prime },๐ค_1,\omega )`$ (50)
$`=\delta _{๐ค,๐ค_1}\delta _{๐ค^{\prime \prime },๐ค_1^{\prime \prime }}\left[f(๐ค)f(๐ค^{\prime \prime })\right]๐_0(๐ค,๐ค^{\prime \prime },\omega ).`$
The related contribution to the current-current correlation function (which is a good approximation for $`\mathrm{\Pi }_{\alpha \alpha }^{\mathrm{intra}}(\omega )`$ at high frequencies) is given by
$`\mathrm{\Pi }_{\alpha \alpha }^{\mathrm{intra},(0)}(\omega )={\displaystyle \frac{1}{\left(\mathrm{}\omega \right)^2}}{\displaystyle \frac{1}{V}}{\displaystyle \underset{\mathrm{๐ค๐ค}^{\prime \prime }\sigma }{}}J_\alpha (๐ค)\left(J_\alpha (๐ค)J_\alpha (๐ค^{\prime \prime })\right)`$
$`\times \left[f(๐ค)f(๐ค^{\prime \prime })\right]|V(๐ค๐ค^{\prime \prime })|^2`$
$`\times {\displaystyle \frac{1}{\mathrm{}}}\left[๐_0(๐ค,๐ค^{\prime \prime },\omega )๐_0(๐ค^{\prime \prime },๐ค,\omega )\right]`$ (51)
$`{\displaystyle \frac{1}{\left(\mathrm{}\omega \right)^2}}{\displaystyle \frac{1}{V}}{\displaystyle \underset{\mathrm{๐ค๐ค}^{\prime \prime }\sigma }{}}J_\alpha (๐ค)\left(J_\alpha (๐ค)J_\alpha (๐ค^{\prime \prime })\right)`$
$`\times \left[f(E(๐ค))f(E(๐ค)+\mathrm{}\omega )\right]|V(๐ค๐ค^{\prime \prime })|^2`$
$`\times {\displaystyle \frac{1}{\mathrm{}}}[๐_0(๐ค,๐ค^{\prime \prime },\omega )+๐_0(๐ค^{\prime \prime },๐ค,\omega )].(47^{})`$
For the impurity scattering processes, the real part of this function is negligible, while the imaginary part is given by
$`\mathrm{Im}\{\mathrm{\Pi }_{\alpha \alpha }^{\mathrm{intra}(0)}(\omega )\}{\displaystyle \frac{1}{\left(\mathrm{}\omega \right)^2}}{\displaystyle \frac{1}{V}}{\displaystyle \underset{๐ค\sigma }{}}J_\alpha ^2(๐ค)`$ (52)
$`\times \left[f(E(๐ค))f(E(๐ค)+\mathrm{}\omega )\right]{\displaystyle \frac{\mathrm{}}{\tau (๐ค,\omega )}},`$
where the electron-hole damping energy is
$`{\displaystyle \frac{\mathrm{}}{\tau (๐ค,\omega )}}={\displaystyle \underset{๐ค^{\prime \prime }}{}}|V(๐ค๐ค^{\prime \prime })|^2\left(1{\displaystyle \frac{J_\alpha (๐ค^{\prime \prime })}{J_\alpha (๐ค)}}\right)`$
$`\times (){\displaystyle \frac{2}{\mathrm{}}}\mathrm{Im}\left\{๐_0(๐ค,๐ค^{\prime \prime },\omega )\right\}.`$ (53)
Furthermore, in this case, the frequency dependent part in $`1/\tau (๐ค,\omega )`$ is negligibly small, and we can put the average over the Fermi surface $`1/\tau (๐ค,0)_{\mathrm{FS}}1/\tau `$ in Eq. (52) instead of $`1/\tau (๐ค,\omega )`$. We finally get KupcicPB2 (15)
$`\mathrm{Im}\{\mathrm{\Pi }_{\alpha \alpha }^{\mathrm{intra},(0)}(\omega )\}{\displaystyle \frac{e^2n_{\mathrm{intra},\alpha }^{\mathrm{eff}}}{m}}{\displaystyle \frac{1}{\omega \tau }},`$ (54)
with
$`n_{\mathrm{intra},\alpha }^{\mathrm{eff}}`$ $`=`$ $`{\displaystyle \frac{m}{e^2}}{\displaystyle \frac{1}{V}}{\displaystyle \underset{๐ค\sigma }{}}J_\alpha ^2(๐ค){\displaystyle \frac{f(๐ค)}{E(๐ค)}}`$ (55)
being the effective number of conduction electrons. The behaviour of $`n_{\mathrm{intra},\mathrm{a}}^{\mathrm{eff}}`$ with band filling $`\delta `$ is shown in Fig. 8 for a few typical values of the ratio $`\mathrm{\Delta }/(2t_a)`$, and compared to the prediction of the free-electron and free-hole approximation. Notice that for $`\mathrm{\Delta }>2t_a2t_b`$ one obtains the well-known result $`n_{\mathrm{intra},\mathrm{a}}^{\mathrm{eff}}\mathrm{sin}k_{\mathrm{F}x}2a=\mathrm{sin}\delta \pi `$ ($`2a`$ is the primitive-cell parameter of the dimerized lattice).
Beyond this approximation the expression (51) requires the evaluation of two coupled integrations over k and $`๐ค^{\prime \prime }`$. In the low-dimensional electronic systems, where the van Hove singularities in the band structure may play important role, this is not a trivial task.
The most outstanding advantage of the present approach is the fact that the effective current vertex $`j_\alpha (๐ค,๐ค^{\prime \prime })`$ consists of two terms, which, together with two terms in $`j_\alpha (๐ค_1^{\prime \prime },๐ค_1)`$, give four contributions to $`\mathrm{\Pi }_{\alpha \alpha }^{\mathrm{intra},(0)}(\omega )`$, two of them representing the so-called self-energy contributions, and the other two the vertex corrections (see Fig. 6(b)) Mahan (7, 25). We shall show next that in the $`n`$th order term in (49) the single-particle vertex corrections and the single-particle self-energy contributions are also treated on equal footing, even if the relaxation processes on impurities are treated beyond the approximation $`1/\tau (๐ค,\omega )1/\tau `$.
#### IV.2.3 Kernel in the low-frequency limit
The effects of the impurity scattering on the correlation function (41) can be exhibited in the following way:
$`๐_2(๐ค,๐ค^{\prime \prime },๐ค_1^{\prime \prime },๐ค_1,\omega )={\displaystyle \underset{๐ช}{}}V(๐ช)[๐_1(๐ค,๐ค^{\prime \prime }๐ช,๐ค_1^{\prime \prime },๐ค_1,\omega )`$
$`๐_1(๐ค+๐ช,๐ค^{\prime \prime },๐ค_1^{\prime \prime },๐ค_1,\omega )].`$ (56)
When combined with this expression, Eqs. (46) and (50) lead to
$`๐_1(๐ค,๐ค^{\prime \prime },๐ค_1^{\prime \prime },๐ค_1,\omega )๐_1^{(0)}(๐ค,๐ค^{\prime \prime },๐ค_1^{\prime \prime },๐ค_1,\omega )`$ (57)
$`+{\displaystyle \frac{1}{\mathrm{}}}๐_0(๐ค,๐ค^{\prime \prime },\omega ){\displaystyle \underset{๐ช}{}}|V(๐ช)|^2`$
$`\times {\displaystyle \frac{1}{\mathrm{}}}\left\{๐_0(๐ค,๐ค^{\prime \prime }+๐ช,\omega )+๐_0(๐ค+๐ช,๐ค^{\prime \prime },\omega )\right\}`$
$`\times \left[๐_1(๐ค,๐ค^{\prime \prime },๐ค_1^{\prime \prime },๐ค_1,\omega )๐_1(๐ค+๐ช,๐ค^{\prime \prime }+๐ช,๐ค_1^{\prime \prime },๐ค_1,\omega )\right].`$
On the right-hand side of this equation, only the self-consistent terms are taken into account. The first term in the brackets represents the single-particle self-energy contributions and the second one the single-particle vertex corrections. It is important to remember again that these corrections are the largest for $`E(๐ค)E(๐ค^{\prime \prime })`$, so that $`๐_0(๐ค,๐ค^{\prime \prime },\omega )`$ can be approximated by $`\omega ^1`$. As a consequence, the solution of this equation can be sought in powers of $`|V(๐ช)|^2/\omega `$. However, due to the dependence on $`๐ค+๐ช`$ and $`๐ค^{\prime \prime }+๐ช`$ of the vertex-corrections term, we have to turn back to the kernel (IV.2.1) and apply several changes to its dummy variables to obtain the self-consistent description of $`_\alpha (๐ค_1^{\prime \prime },๐ค_1,\omega )`$ and the desired recurrence relations between the contributions $`๐_1^{(2n)}(๐ค,๐ค^{\prime \prime },๐ค_1^{\prime \prime },๐ค_1,\omega )`$. The kernel is described by
$`_\alpha (๐ค_1^{\prime \prime },๐ค_1,\omega )_\alpha ^{(0)}(๐ค_1^{\prime \prime },๐ค_1,\omega )`$
$`={\displaystyle \frac{1}{\omega }}{\displaystyle \underset{\mathrm{๐ค๐ค}^{\prime \prime }}{}}\left|V(๐ค๐ค^{\prime \prime })\right|^2\left[J_\alpha (๐ค^{\prime \prime })J_\alpha (๐ค)\right]`$
$`\times ๐_1(๐ค,๐ค^{\prime \prime },๐ค_1^{\prime \prime },๐ค_1,\omega )\mathrm{\Sigma }(๐ค,๐ค^{\prime \prime },\omega ),`$ (58)
where
$`\mathrm{}\mathrm{\Sigma }(๐ค,๐ค^{\prime \prime },\omega )`$ (59)
$`={\displaystyle \underset{๐ช}{}}|V(๐ช๐ค^{\prime \prime })|^2\left(1{\displaystyle \frac{J_\alpha (๐ช)}{J_\alpha (๐ค^{\prime \prime })}}\right){\displaystyle \frac{1}{\mathrm{}}}๐_0(๐ค,๐ช,\omega )`$
$`{\displaystyle \underset{๐ช}{}}|V(๐ช๐ค)|^2\left(1{\displaystyle \frac{J_\alpha (๐ช)}{J_\alpha (๐ค)}}\right){\displaystyle \frac{1}{\mathrm{}}}๐_0(๐ช,๐ค^{\prime \prime },\omega )`$
is the electron-hole self-energy and $`_\alpha ^{(0)}(๐ค_1^{\prime \prime },๐ค_1,\omega )`$ is the kernel corresponding to the replacement of $`๐_1(๐ค,๐ค^{\prime \prime },๐ค_1^{\prime \prime },๐ค_1,\omega )`$ in Eq. (IV.2.1) by $`๐_1^{(0)}(๐ค,๐ค^{\prime \prime },๐ค_1^{\prime \prime },๐ค_1,\omega )`$. The only approximation made in the derivation of Eq. (58) is
$`J_\alpha (๐ค๐ค^{\prime \prime }+๐ช)J_\alpha (๐ช)\left[J_\alpha (๐ค)J_\alpha (๐ค^{\prime \prime })\right]{\displaystyle \frac{J_\alpha (๐ช)}{J_\alpha (๐ค^{\prime \prime })}},`$
(60)
which allows a simple description of the vertex corrections in $`_\alpha (๐ค_1^{\prime \prime },๐ค_1,\omega )`$, and which treats correctly the disappearance of the forward scattering contributions ($`๐ค๐ค^{\prime \prime }`$) to both $`\mathrm{\Sigma }(๐ค,๐ค^{\prime \prime },\omega )`$ and $`_\alpha (๐ค_1^{\prime \prime },๐ค_1,\omega )`$. As briefly discussed at the end of Sec. IV B 5, this approximation is directly related to the restrictions enforced by the the continuity equation for the charge density. Similarly, the definition (49), together with the self-consistent relations (57) and (58), gives rise to the recurrence relations illustrated in Fig. 7(d):
$`๐_1^{(2n)}(๐ค,๐ค^{\prime \prime },๐ค_1^{\prime \prime },๐ค_1,\omega )`$ (61)
$`={\displaystyle \frac{\mathrm{\Sigma }(๐ค,๐ค^{\prime \prime },\omega )}{\omega }}๐_1^{(2n2)}(๐ค,๐ค^{\prime \prime },๐ค_1^{\prime \prime },๐ค_1,\omega )`$
$`=\left({\displaystyle \frac{\mathrm{\Sigma }(๐ค,๐ค^{\prime \prime },\omega )}{\omega }}\right)^n๐_1^{(0)}(๐ค,๐ค^{\prime \prime },๐ค_1^{\prime \prime },๐ค_1,\omega ).`$
#### IV.2.4 Memory-function approximation
The simplest way to solve Eq. (58) is to replace $`\mathrm{\Sigma }(๐ค,๐ค^{\prime \prime },\omega )`$ by its imaginary part averaged over the Fermi surface, $`\mathrm{i}/\tau (\omega )`$. As mentioned above, for the impurity scattering processes, the real part of $`\mathrm{\Sigma }(๐ค,๐ค^{\prime \prime },\omega )`$ can be ignored. Even if $`\mathrm{Re}\{\mathrm{\Sigma }(๐ค,๐ค^{\prime \prime },\omega )\}`$ is not small, we can turn back to the electronic Hamiltonian and tray to include the $`\mathrm{Re}\{\mathrm{\Sigma }(๐ค,๐ค^{\prime \prime },\omega )\}`$ effects into the effective single-particle Hamiltonian, and to diagonalize it, as we did here with the scattering processes on the dimerization potential $`H_0^{}`$. The real part of the new self-energy $`\mathrm{\Sigma }(๐ค,๐ค^{\prime \prime },\omega )`$ is minimized in this way, with only the imaginary part playing an important role in Eq. (61).
In this case, we obtain the expression
$`_\alpha (๐ค_1^{\prime \prime },๐ค_1,\omega )`$ $``$ $`_\alpha ^{(0)}(๐ค_1^{\prime \prime },๐ค_1,\omega ){\displaystyle \frac{\omega }{\omega +\mathrm{i}/\tau (\omega )}},`$ (62)
which leads to the well-known results of the memory-function approximation KupcicPB2 (15, 25, 26)
$`\mathrm{\Pi }_{\alpha \alpha }^{\mathrm{intra}}(\omega )`$ $``$ $`{\displaystyle \frac{e^2n_{\mathrm{intra},\alpha }^{\mathrm{eff}}}{m}}{\displaystyle \frac{\mathrm{i}/\tau (\omega )}{\omega +\mathrm{i}/\tau (\omega )}},`$
$`\sigma _{\alpha \alpha }^{\mathrm{intra}}(\omega )`$ $``$ $`{\displaystyle \frac{\mathrm{i}}{\omega }}{\displaystyle \frac{e^2n_{\mathrm{intra},\alpha }^{\mathrm{eff}}}{m}}{\displaystyle \frac{\omega }{\omega +\mathrm{i}/\tau (\omega )}},`$ (63)
with $`\mathrm{}/\tau (\omega )`$ being the intraband memory (relaxation) function. This result is consistent with the causality requirement,
$`\mathrm{Re}\{\sigma _{\alpha \alpha }^{\mathrm{intra}}(\omega )\}`$ $`=`$ $`\mathrm{Re}\{\sigma _{\alpha \alpha }^{\mathrm{intra}}(\omega )\},`$
$`\mathrm{Im}\{\sigma _{\alpha \alpha }^{\mathrm{intra}}(\omega )\}`$ $`=`$ $`\mathrm{Im}\{\sigma _{\alpha \alpha }^{\mathrm{intra}}(\omega )\},`$ (64)
provided that $`\tau (\omega )=\tau `$. If $`\tau (\omega )`$ is frequency dependent, the corresponding $`\mathrm{Re}\{\mathrm{\Sigma }(๐ค,๐ค^{\prime \prime },\omega )\}`$ is non-zero, but the result (63) is still acceptable. For $`\mathrm{Re}\{\mathrm{\Sigma }(๐ค,๐ค^{\prime \prime },\omega )\}`$ not too large we can introduce the effects of $`\mathrm{Re}\{\mathrm{\Sigma }(๐ค,๐ค^{\prime \prime },\omega )\}`$ through the mass redefinition $`mm(\omega )`$ through the KramersโKronig relations. The result is the generalized Drude formula for the intraband optical conductivity.
Here we show two important results. First, the memory-function approximation, which in the traditional form has not been found to be transparent, can be understood as a simple replacement $`\mathrm{\Sigma }(๐ค,\omega )\mathrm{i}/\tau (\omega )`$ of the exact-summation result given below. Second, the memory-function results are acceptable even in the cases with the pronounced optical excitations across a gap (or pseudogap) (where $`n_{\mathrm{intra},\alpha }^{\mathrm{eff}}n`$; see, for example, the hole-like semiconducting regime in Fig. 8 at $`\delta 1`$), provided that the electron group velocity $`v_\alpha (๐ค)=J_\alpha (๐ค)/e`$ in Eq. (55) is determined using the relation (77). The intraband conductivity spectrum obtained in this way is related with the interband conductivity spectrum through the well-controlled conductivity sum rule KupcicPB2 (15).
#### IV.2.5 Exact summation
In the case where the wave vector dependence of the imaginary part of $`\mathrm{\Sigma }(๐ค,๐ค^{\prime \prime },\omega )`$ is significant and the real part of $`\mathrm{\Sigma }(๐ค,๐ค^{\prime \prime },\omega )`$ is not too large the full recurrence relations for $`๐_1^{(2n)}(๐ค,๐ค^{\prime \prime },๐ค_1^{\prime \prime },๐ค_1,\omega )`$ can be used to obtain the intraband current-current correlation function. The result is
$`\mathrm{\Pi }_{\alpha \alpha }^{\mathrm{intra}}(\omega )={\displaystyle \frac{1}{V}}{\displaystyle \underset{\mathrm{๐ค๐ค}^{\prime \prime }\sigma }{}}\left[J_\alpha (๐ค)J_\alpha (๐ค^{\prime \prime })\right]^2{\displaystyle \frac{|V(๐ค๐ค^{\prime \prime })|^2}{\mathrm{}\omega }}`$
$`\times \left[f(๐ค)f(๐ค^{\prime \prime })\right]{\displaystyle \frac{1}{\mathrm{}}}{\displaystyle \frac{๐_0(๐ค,๐ค^{\prime \prime },\omega )}{\mathrm{}\omega +\mathrm{}\mathrm{\Sigma }(๐ค,๐ค^{\prime \prime },\omega )}},`$ (65)
$`{\displaystyle \frac{1}{V}}{\displaystyle \underset{\mathrm{๐ค๐ค}^{\prime \prime }\sigma }{}}J_\alpha ^2(๐ค){\displaystyle \frac{f(๐ค)}{E(๐ค)}}{\displaystyle \frac{1}{\mathrm{}\omega +\mathrm{}\mathrm{\Sigma }(๐ค,๐ค^{\prime \prime },\omega )}}`$
$`\times |V(๐ค๐ค^{\prime \prime })|^2\left(1{\displaystyle \frac{J_\alpha (๐ค^{\prime \prime })}{J_\alpha (๐ค)}}\right)`$
$`\times {\displaystyle \frac{1}{\mathrm{}}}[๐_0(๐ค,๐ค^{\prime \prime },\omega )+๐_0(๐ค^{\prime \prime },๐ค,\omega )].(61^{})`$
Physically the most important case corresponds to the approximation $`\mathrm{\Sigma }(๐ค,๐ค^{\prime \prime },\omega )\mathrm{\Sigma }(๐ค,๐ค,\omega )=\mathrm{\Sigma }(๐ค,\omega )`$. In this case, we obtain
$`\mathrm{\Pi }_{\alpha \alpha }^{\mathrm{intra}}(\omega ){\displaystyle \frac{1}{V}}{\displaystyle \underset{๐ค\sigma }{}}J_\alpha ^2(๐ค){\displaystyle \frac{f(๐ค)}{E(๐ค)}}{\displaystyle \frac{\mathrm{\Sigma }(๐ค,\omega )}{\omega +\mathrm{\Sigma }(๐ค,\omega )}}.(61^{\prime \prime })`$
The final form of the optical conductivity comes from Eqs. ($`61^{\prime \prime }`$) and (11)
$`\sigma _{\alpha \alpha }^{\mathrm{intra}}(\omega ){\displaystyle \frac{\mathrm{i}}{\omega }}{\displaystyle \frac{1}{V}}{\displaystyle \underset{๐ค\sigma }{}}J_\alpha ^2(๐ค)(){\displaystyle \frac{f(๐ค)}{E(๐ค)}}{\displaystyle \frac{\omega }{\omega +\mathrm{\Sigma }(๐ค,\omega )}},`$
(66)
which is the result identical to the result of the longitudinal response theory.
The longitudinal response theory gives a simpler way to obtain the same $`\sigma _{\alpha \alpha }^{\mathrm{intra}}(\omega )`$. The difference between the two approaches is in the way how the continuity equation connecting the charge density and current density fluctuations is treated. The consideration of the direct processes of the wave vector q in the longitudinal approach allows a more precise treatment of the continuity equation, in the way analogous to the Landau response theory Pines (11, 12). But here only an approximate solution is possible, since the theory is formulated in terms of the indirect intraband processes (of the wave vector $`๐ค๐ค^{\prime \prime }`$).
#### IV.2.6 Zero-frequency limit
First significant consequence of the present exact summation method is the behaviour of the intraband optical conductivity in the zero-frequency limit. When the real part of $`\mathrm{\Sigma }(๐ค,\omega )`$ is small enough, we can write
$`\mathrm{\Sigma }(๐ค,\omega )`$ $``$ $`\mathrm{i}\mathrm{\Sigma }^{\prime \prime }(๐ค,0)\mathrm{i}/\tau (๐ค),`$ (67)
resulting in the DC conductivity which is equal to the well-known Boltzmann result Mahan (7, 10)
$`\sigma _{\alpha \alpha }^{\mathrm{intra}}(0)\sigma _{\alpha \alpha }^{\mathrm{DC}}`$ $`=`$ $`(){\displaystyle \frac{1}{V}}{\displaystyle \underset{๐ค\sigma }{}}J_\alpha ^2(๐ค){\displaystyle \frac{f(๐ค)}{E(๐ค)}}\tau (๐ค)`$ (68)
$`=`$ $`{\displaystyle \frac{e^2\tau _0}{mV_0}}{\displaystyle \frac{m}{m_{aa}}}\stackrel{~}{n}_{\mathrm{intra},\alpha }^{\mathrm{eff}}.`$
Here
$`\stackrel{~}{n}_{\mathrm{intra},\alpha }^{\mathrm{eff}}`$ $`=`$ $`{\displaystyle \frac{m_{aa}}{m}}V_0n_{\mathrm{intra},\alpha }^{\mathrm{eff}}`$
$`=`$ $`(){\displaystyle \frac{m_{aa}}{N}}{\displaystyle \underset{๐ค\sigma }{}}v_\alpha ^2(๐ค){\displaystyle \frac{f(๐ค)}{E(๐ค)}}{\displaystyle \frac{\tau (๐ค)}{\tau _0}},(51^{})`$
is the effective number of conduction electrons shown in the dimensionless form and $`\tau _0=\tau (๐ค=0)`$ is the temperature dependent $`๐ค=0`$ relaxation time. The relation (66), together with Eq. (70), gives the complete description of the optical conductivity in a general multiband model, with the symmetry of the intra- and interband current vertices playing an essential role. Thus, Eq. (68) provides the direct link between the low-frequency optical conductivity and various transport coefficients not only in the single-band but also in the multiband models.
### IV.3 Interband optical conductivity
The approximation in which the $`๐_2(๐ค,๐ค^{\prime \prime },๐ค_1,\omega )`$ term in Eq. (35) is omitted leads to the ideal interband current-current correlation function and to the ideal interband conductivity characterized by a sharp threshold at the energy $`E_{\underset{ยฏ}{C}}(๐ค_\mathrm{F})E_C(๐ค_\mathrm{F})`$ LRA (1, 2, 3, 15). The former one is given by
$`\mathrm{\Pi }_{\alpha \alpha }^{\mathrm{inter}}(\omega )`$ $`=`$ $`{\displaystyle \frac{1}{V}}{\displaystyle \underset{๐ค\sigma }{}}|J_\alpha ^{\underset{ยฏ}{C}C}(๐ค)|^2\{{\displaystyle \frac{f_C(๐ค)f_{\underset{ยฏ}{C}}(๐ค)}{\mathrm{}\omega E_{\underset{ยฏ}{C}C}(๐ค)+\mathrm{i}\mathrm{}\eta }}`$ (69)
$`+{\displaystyle \frac{f_{\underset{ยฏ}{C}}(๐ค)f_C(๐ค)}{\mathrm{}\omega +E_{\underset{ยฏ}{C}C}(๐ค)+\mathrm{i}\mathrm{}\eta }}\},`$
with $`E_{\underset{ยฏ}{C}C}(๐ค)=E_{\underset{ยฏ}{C}}(๐ค)E_C(๐ค)`$. The latter one comes from Eqs. (11) and (69)
$`\sigma _{\alpha \alpha }^{\mathrm{inter}}(\omega )`$ $`=`$ $`{\displaystyle \frac{\mathrm{i}}{\omega }}{\displaystyle \frac{1}{V}}{\displaystyle \underset{๐ค\sigma }{}}{\displaystyle \frac{\mathrm{}\omega |J_\alpha ^{\underset{ยฏ}{C}C}(๐ค)|^2}{E_{\underset{ยฏ}{C}C}(๐ค)}}`$ (70)
$`\times {\displaystyle \frac{2\left[f_C(๐ค)f_{\underset{ยฏ}{C}}(๐ค)\right]}{\mathrm{}\omega +\mathrm{i}\mathrm{}\eta E_{\underset{ยฏ}{C}C}^2(๐ค)/\left(\mathrm{}\omega \right)}},`$
with $`\eta =0`$ put in the numerator and $`\eta 0^+`$ in the denominator. The leading term in the interband electron-hole self-energy
$`\mathrm{}\mathrm{\Sigma }_{LL^{}}(๐ค,\omega )`$ $`=`$ $`{\displaystyle \underset{๐ช}{}}[|V^{L^{}L^{}}(๐ช๐ค)|^2{\displaystyle \frac{1}{\mathrm{}}}๐_0^{LL^{}}(๐ค,๐ช,\omega )`$ (71)
$`+|V^{LL}(๐ช๐ค)|^2{\displaystyle \frac{1}{\mathrm{}}}๐_0^{LL^{}}(๐ช,๐ค,\omega )]`$
comes from the single-particle self-energy contributions. A reasonable generalization for the interband optical conductivity is given by Eq. (70) with the replacement $`\eta `$ by $`\mathrm{Im}\{\mathrm{\Sigma }_{\underset{ยฏ}{C}C}(๐ค,\omega )\}`$.
The interplay between the self-energy and vertex-corrections terms in $`\mathrm{\Sigma }_{\underset{ยฏ}{C}C}(๐ค,\omega )`$, the correct treatment of the indirect interband optical excitations, as well as the role of the effective mass theorem in resolving all these issues will be explained elsewhere KupcicUP (12). It should be noticed here that the Landau-like function (70) gives the in-gap optical conductivity slightly different from the corresponding Lindhard-like function, as easily seen by comparing Fig. 9 with Fig. 4 reported in Ref. KupcicPB2 (15). The understanding of the difference between these two results is of general importance as well, and will be discussed in detail in Ref. KupcicUP (12).
## V Comparison with experiments
In the simples limit, $`\mathrm{\Sigma }_{CC}(๐ค,\omega )\mathrm{i}\mathrm{\Gamma }_{\mathrm{intra}}`$ and $`\mathrm{\Sigma }_{C\underset{ยฏ}{C}}(๐ค,\omega )\mathrm{i}\mathrm{\Gamma }_{\mathrm{inter}}`$, the optical conductivity of the present two-band model,
$`\sigma _{\alpha \alpha }(\omega )`$ $`=`$ $`\sigma _{\alpha \alpha }^{\mathrm{intra}}(\omega )+\sigma _{\alpha \alpha }^{\mathrm{inter}}(\omega ),`$ (72)
is a function of the Fermi wave vector $`๐ค_\mathrm{F}`$, three band parameters, $`t_a`$, $`t_b`$ and $`\mathrm{\Delta }`$, and two damping energies, $`\mathrm{}\mathrm{\Gamma }_{\mathrm{intra}}`$ and $`\mathrm{}\mathrm{\Gamma }_{\mathrm{inter}}`$. For $`k_{\mathrm{F}x}0.5\pi /a`$, the model illustrates optical properties of various Q1D imperfectly nested CDW systems (including both the ordered CDW state and the pseudogap effects at temperatures above the critical temperature $`T_{\mathrm{CDW}}`$). In this section, we shall briefly discuss a few qualitative results important to the Q1D CDW systems. First, the temperature dependence of the DC and optical conductivity in the ordered CDW state is discussed for the strictly 1D case ($`t_b0`$). Then we contrast the interband conductivity found here to the oversimplified semiconducting optical conductivity usually encountered in the textbooks LRA (1, 2, 3).
### V.1 DC conductivity in the ordered CDW state
The temperature dependence in $`k_\mathrm{B}T`$, $`\mathrm{\Delta }(T)`$, $`\mathrm{\Gamma }_{\mathrm{intra}}(T)`$ and $`\mathrm{\Gamma }_{\mathrm{inter}}(T)`$ is responsible for the transfer with increasing/decreasing temperature of the optical conductivity spectra between the intraband and interband channels. This effect is particularly large in the vicinity of the metal-to-insulator phase transition. It should be noticed that, in the approximation in which $`\tau (๐ค)\tau _0=1/\mathrm{\Gamma }_{\mathrm{intra}}`$ is set in Eq. ($`51^{}`$), the temperature dependence of $`\sigma _{\alpha \alpha }^{\mathrm{DC}}`$ is given by the product of the effective number of conduction electrons $`\stackrel{~}{n}_{\mathrm{intra},\alpha }^{\mathrm{eff}}(T)`$ and the relaxation time $`\tau _0(T)`$. We have two adjustable parameters at any temperature, and the analysis of the DC conductivity data is thus possible only by the combination with the optical conductivity measurements. The latter method allows also the determination of the magnitude of CDW order parameter $`\mathrm{\Delta }_0`$ and the critical exponent $`\beta `$ in $`\mathrm{\Delta }(T)=\mathrm{\Delta }_0\left(1T/T_{\mathrm{CDW}}\right)^\beta `$, as well as the bond energy $`t_a`$ and the damping energy $`\mathrm{\Gamma }_{\mathrm{inter}}`$.
In Figs. 9 and 10 the optical conductivity normalized to the DC conductivity at $`\mathrm{\Delta }=0`$ and the temperature dependence of the DC conductivity are shown for typical values of the parameters.
### V.2 Optical conductivity of a simple semiconductor
The typical result for the interband conductivity in the ordered CDW state is shown in Fig. 11 (solid and dotted curves), and compared to the data measured in the blue bronze K<sub>0.3</sub>MoO<sub>3</sub> (diamonds) Degiorgi (5). Notice that the gauge-invariance factor $`\mathrm{}\omega /E_{\underset{ยฏ}{C}C}(๐ค)`$ in Eq. (70) ensures the disappearance of $`\sigma _{\alpha \alpha }^{\mathrm{inter}}(0)`$ at $`T0`$, independently of the value of the phenomenologically introduced damping energy $`\mathrm{}\mathrm{\Gamma }_{\mathrm{inter}}`$. The dashed curve is the prediction of the usual optical model for semiconductors LRA (1, 2, 3), with the factor $`\mathrm{}\omega /E_{\underset{ยฏ}{C}C}(๐ค)`$ absent, which is characterized by a significant (but non-physical) contribution to $`\sigma _{\alpha \alpha }^{\mathrm{DC}}`$ at $`T0`$.
## VI Conclusion
In this article, the response of the conduction electrons in a Q1D multiband model to the transverse electromagnetic fields has been studied in presence of two scattering mechanisms. In contrast to the coherent scattering on the site-energy-dimerization potential, the impurity scattering processes give negligibly small contribution to the real part of the electron self-energy, but dominate in the relaxation processes in both the intra- and interband optical excitations. We determine the current-current correlation function, i.e. the optical conductivity of the related two-band model as a function of band filling. The transverse equation of motion approach has been used to collect the most singular contributions to the optical conductivity. It is shown that the present multiband optical analysis represents a generalization of the usual force-force correlation function method, and that in the DC limit it approaches correctly the results of Boltzmann equations, due to its gauge-invariant form. The present optical conductivity model gives the frequency and temperature dependences of the single-particle contributions to the optical conductivity spectra in the ordered CDW state which are consistent with both the experimental observation and the prediction of the longitudinal response theory. It is explained that, while the response to the longitudinal fields is associated with the direct electron-hole pair excitations, the response to the transverse electromagnetic fields can be understood in terms of the indirect (impurity-assisted) electron-hole pair excitations.
## Acknowledgements
This research was supported by the Croatian Ministry of Science and Technology under the project 0119-256. One of the authors (I.K.) would like to acknowledge the hospitality of the Institute of Physics of Complex Matter, Lausanne, where parts of this work were completed.
## Appendix A Current and static Raman vertices
The vertex functions in the expression (29) depend on the unperturbed vertices $`J_\alpha ^{ll}(๐ค)=(e/\mathrm{})\epsilon _l(๐ค)/k_\alpha `$ and $`\gamma _{\alpha \alpha }^{ll}(๐ค;2)=(m/\mathrm{}^2)^2\epsilon _l(๐ค)/k_\alpha ^2`$ in the way
$`J_\alpha ^{LL^{}}(๐ค)`$ $`=`$ $`{\displaystyle \underset{l}{}}U_๐ค(l,L)U_๐ค^{}(l,L^{})J_\alpha ^{ll}(๐ค),`$
$`\gamma _{\alpha \alpha }^{LL^{}}(๐ค;2)`$ $`=`$ $`{\displaystyle \underset{l}{}}U_๐ค(l,L)U_๐ค^{}(l,L^{})\gamma _{\alpha \alpha }^{ll}(๐ค;2).`$ (73)
For the $`\alpha =a`$ polarization of the electromagnetic fields, the result is
$`J_a^{\underset{ยฏ}{C}\underset{ยฏ}{C},CC}(๐ค)`$ $`=`$ $`\left[u^2(๐ค)v^2(๐ค)\right]J_a^{cc}(๐ค)`$ (74)
$`=`$ $`\mathrm{cos}\phi (๐ค)J_a^{cc}(๐ค),`$
$`J_a^{\underset{ยฏ}{C}C}(๐ค)`$ $`=`$ $`2u(๐ค)v(๐ค)J_a^{cc}(๐ค)`$
$`=`$ $`\mathrm{sin}\phi (๐ค)J_a^{cc}(๐ค),(\mathrm{A2}^{})`$
$`\gamma _{aa}^{\underset{ยฏ}{C}\underset{ยฏ}{C},CC}(๐ค;2)`$ $`=`$ $`\left[u^2(๐ค)v^2(๐ค)\right]\gamma _{aa}^{cc}(๐ค;2)`$ (75)
$`=`$ $`\mathrm{cos}\phi (๐ค)\gamma _{aa}^{cc}(๐ค;2).`$
Similarly, for $`\alpha =b`$
$`J_b^{\underset{ยฏ}{C}\underset{ยฏ}{C},CC}(๐ค)=J_b^{cc}(๐ค),`$ $`J_b^{\underset{ยฏ}{C}C}(๐ค)=0,`$
$`\gamma _{bb}^{\underset{ยฏ}{C}\underset{ยฏ}{C},CC}(๐ค;2)`$ $`=`$ $`\gamma _{bb}^{cc}(๐ค;2).`$ (76)
Finally, using Eqs. (18), (74) and (76) we can check the Ward identity Mahan (7) which connects the intraband current vertex $`J_\alpha ^{LL}(๐ค)`$ with the electron group velocity $`v_\alpha ^L(๐ค)`$
$`J_\alpha ^{LL}(๐ค)`$ $`=`$ $`ev_\alpha ^L(๐ค)={\displaystyle \frac{e}{\mathrm{}}}{\displaystyle \frac{E_L(๐ค)}{k_\alpha }}.`$ (77)
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# 1 Introduction
## 1 Introduction
Various arguments suggest that the description of black holes should be an important aspect of a quantum-gravity theory, and that some key operatively-meaningful (not merely formal) differences between alternative theories should emerge as we establish, within each approach, how the singularity, the evaporation, the โthermodynamicsโ, and the โinformation paradoxโ are handled. A similar role could be played, as stressed in the recent literature , by the analysis of the energy-momentum dispersion relation and the position-momentum uncertainty principle. Different approaches to the quantum-gravity problem lead to different expectations for what concerns the possibility of a MDR (a modified energy-momentum dispersion relation) and the possibility of a GUP (a generalized position-momentum uncertainty principle). In particular, in the study of Loop Quantum Gravity and of models based on noncommutative geometry there has been strong interest in some candidate modifications of the energy-momentum dispersion relation. Generalized uncertainty principles have been considered primarily in the literature on String Theory and on models based on noncommutative geometry . The form of the energy-momentum dispersion relation and of the position-momentum uncertainty relation can therefore be used to characterize alternative approaches to the quantum-gravity problem.
Two of us were recently involved in research exploring a possible link between the predictions that a quantum-gravity theory makes for black-hole thermodynamics and the predictions that the same theory makes for the energy-momentum dispersion relation and the position-momentum uncertainty relation. By establishing the nature of such a link one would, in our opinion, obtain a valuable characterization of the type of internal logical consistency that various aspects of a quantum-gravity theory should satisfy. The study reported in Ref. was not the first to explore the possible role of MDRs and/or GUPs for black holes, as we discuss in greater detail here in Section 9 where we comment of some relevant references , but all of these studies, including Ref. , focused on one or another aspect of the possible interplay between MDR/GUP results and black holes, without attempting to obtain a wider picture. We here work in the spirit of Ref. , but we attempt to give the first elements of a general analysis of some key characteristics of black-hole physics, as affected by some scenarios for a MDR or a GUP.
Sections 2, 3 and 4, set the stage, by reviewing some results in the MDR/GUP literature and revisiting the point already made in Ref. , which concerns an apparent link between the log-area terms in the entropy-area relation for black holes and certain formulations of the MDR and the GUP. In Section 5 we explore, still within the working assumptions adopted in Ref. , the implications of a MDR and/or a GUP for the Bekenstein entropy bound and for the Generalized Second Law of thermodynamics. We find that the implications are significant and we conjecture that they should also not be negligible in the analysis of other entropy-bound proposals. Section 6 considers a role for MDR/GUP modifications in the analysis of the black-body radiation spectrum, and again exposes some significant changes with respect to the standard picture, including the possibility that the characteristic frequency of black-body radiation at given temperature $`T`$ might have a dependence on $`T`$ such that in the infinite-temperature limit the characteristic frequency would take a finite (Planckian) value. It is then perhaps not surprising that in the analysis of the black-hole evaporation process, discussed in Section 7, we also find some characteristic MDR/GUP-induced new features, such as the possibility that the energy flux emitted by the black hole might diverge when the black-hole mass reaches a certain finite (Planckian) value. In Section 8, we comment on one key aspect which might deserve further consideration: for these theories with MDRs and/or GUPs there has been some speculation that the speed of massless particles might be different from the familiar speed-of-light scale value of $`c`$. In Sections 1-7 we assume throughout that $`c`$ still is the speed of massless particles, but in Section 8 we establish how the analysis of black-body radiation would be changed if one implemented some alternatives considered in the literature. In Section 9 we compare our analysis with other studies which have considered the implications of a MDR or a GUP for some aspects of black-hole physics. Section 10 concludes the paper with some remarks on the outlook of this research programme.
## 2 MDRs and GUPs in Quantum Gravity and implications for a Planck-scale particle-localization limit
### 2.1 MDRs and GUPs in Quantum Gravity
In the study of the Quantum-Gravity problem the emergence of modified energy-momentum relations and/or generalized position-momentum uncertainty principles, although of course not guaranteed, can be motivated on general grounds, and also finds support in the direct analysis of certain Quantum-Gravity scenarios.
The hypothesis of modified energy-momentum dispersion relations is understandably popular among those adopting a โspacetime foamโ intuition in the study of the quantum-gravity problem, especially when an analogy between spacetime foam and some more familiar forms of medium (such as certain crystal structures of interest in condensed-matter studies) is proposed. It is then expected that wave dispersion โin vacuoโ (in the spacetime foam) might resemble wave dispersion in other media. A modified dispersion relation can also be favoured by the expectation, shared by many researchers of the field, that the Planck length might fundamentally set the minimum allowed value for wavelengths. A nonlinear relationship between energy and (space-) momentum can be easily adjusted in such a way that in the infinite-energy limit the momentum saturates to the Planck-scale value (and wavelength saturates to the Planck-length value). This possibility has become more attractive with the recent realization that a modified energy-momentum dispersion relation can also be introduced as an observer-independent law<sup>1</sup><sup>1</sup>1But this usually requires introducing a nonlinear deformation of the action of Lorentz boosts., in which case the Planckian minimum-wavelength hypothesis can be introduced as a physical law valid in every frame. The analysis of some quantum-gravity scenarios, even in cases in which the emergence of modified energy-momentum relations was not intended in the original setup of the framework, has shown some explicit mechanisms for the emergence of modified dispersion relations. This is particularly true of some approaches based on noncommutative geometry and within the Loop-Quantum-Gravity approach . In most cases one is led to consider a dispersion relation of the type<sup>2</sup><sup>2</sup>2We denote with $`m`$, as conventional, the rest energy of the particle. The mass parameter $`\mu `$ on the right-hand side is directly related to the rest energy, but $`\mu m`$ if the $`\alpha _i`$ do not all vanish. For example, if $`\alpha _10`$ but $`\alpha _i=0`$ for every $`i2`$ one of course obtains $`\mu ^2=m^2+\alpha _1L_pm^3`$. This needed to be clarified since it is relevant for more general analyses of MDRs, but in our study we are always concerned with particles which are either massless or anyway are analyzed at energies such that the mass can be neglected, and therefore both $`\mu `$ and $`m`$ will never actually enter our key formulas.
$$\stackrel{}{p}^2=f(E,m;L_p)E^2\mu ^2+\alpha _1L_pE^3+\alpha _2L_p^2E^4+O\left(L_p^3E^5\right),$$
(1)
where $`f`$ is the function that gives the exact dispersion relation, and on the right-hand side we just assumed the applicability of a Taylor-series expansion for $`E1/L_p`$. The coefficients $`\alpha _i`$ can take different values in different Quantum-Gravity proposals.
The fact that these Planck-scale-deformed dispersion relations may have observably large consequences in some (however rare) physical contexts has led to interest in this research also from the perspective of phenomenology .
The situation concerning the possibility of a generalized position-momentum uncertainty principle is rather similar. On general grounds it can be motivated by the intuition that the solution of the quantum-gravity problem might require the introduction of an absolute Planckian limit on the size of the collision region, applicable to high-energy microscopic collision processes. For example, a GUP of the form
$$\delta x\frac{1}{\delta p}+\alpha L_p^2\delta p+O(L_p^3\delta p^2),$$
(2)
which has been derived within the String Theory approach to the quantum-gravity problem , is such that at small $`\delta p`$ one finds the standard dependence of $`\delta x`$ on $`\delta p`$ ($`\delta x`$ gets smaller as $`\delta p`$ increases) but for large $`\delta p`$ the Planckian correction term becomes significant and keeps $`\delta xL_p`$. Within String Theory the coefficient $`\alpha `$ should take a value of roughly the ratio between the square of the string length and the square of the Planck length, but this of course might work out differently in other Quantum-Gravity proposals.
While in the parametrization of (1) we included a possible correction term suppressed only by one power of the Planck length, in (2) such a linear-in-$`L_p`$ is assumed not to be present. This reflects the status of the presently-available literature: for the MDR a large number of alternative formulations, including some with the linear-in-$`L_p`$ term, are being considered, as they find support in different approaches to the quantum-gravity problem (and different preliminary results adopting alternative approximation schemes within a given approach), whereas all the discussions of a GUP assume that the leading-order correction should be proportional to the square of $`L_p`$.
### 2.2 MDR and a Planck-scale particle-localization limit
The analysis reported in Ref. exposed a previously unnoticed common feature of MDR and GUP scenarios. This has to do with a Planck-scale limit on the localization of a particle, and an associated modification of the Bekenstein argument for a area-entropy black-hole relation.
Arguably the closest starting point for the construction of the correct Quantum Gravity should be Quantum Field Theory, and within Quantum Field Theory the most striking quantum effect concerns an absolute limit on the localization of a particle of energy $`E`$, codified in the relation $`E\frac{1}{\delta x}`$. While in nonrelativistic quantum mechanics a particle of any energy can always be sharply localized (at the price of renouncing to all information on the conjugate momentum), within Quantum Field Theory only in the infinite-energy limit a particle can be sharply localized. And among those studying the quantum-gravity problem one frequently encounters the intuition that at the Quantum-Gravity level the idealization of sharp localization should disappear completely.
In the spirit of Ref. one can attempt to codify this quantum-gravity intuition in a relation of the type
$$E\frac{1}{\delta x}\left(1\mathrm{\Delta }(L_p,\delta x)\right)$$
(3)
where $`\mathrm{\Delta }`$ is some function of $`L_p`$ and $`\delta x`$, perhaps such that $`E\mathrm{}`$ already at some finite value of $`\delta x`$ (so that the idealization $`\delta x0`$ is excluded). And it was observed in Ref. that both the idea of a MDR and the idea of a GUP would support a formula of the type (3), with nonzero $`\mathrm{\Delta }`$.
Let us briefly review this analysis reported in Ref. , starting with the case of a MDR of the type (1). We can follow the familiar derivation of the relation $`E\frac{1}{\delta x}`$, substituting, where necessary, the standard special-relativistic dispersion relation with its Planck-scale modified version. It is convenient to start by focusing on the case of a particle of mass $`M`$ at rest, whose position is being measured by a procedure involving a collision with a photon of energy $`E_\gamma `$ and momentum $`p_\gamma `$. According to Heisenbergโs uncertainty principle, in order to measure the particle position with precision $`\delta x`$ one should use a photon with momentum uncertainty $`\delta p_\gamma \frac{1}{\delta x}`$. Following the standard argument , one takes this $`\delta p_\gamma \frac{1}{\delta x}`$ relation and converts it into the relation $`\delta E_\gamma \frac{1}{\delta x}`$ using the special relativistic dispersion relation. Finally $`\delta E_\gamma \frac{1}{\delta x}`$ is converted into the relation $`M\frac{1}{\delta x}`$ because the measurement procedure requires $`\delta EM`$, in order to ensure that the relevant energy uncertainties are not large enough to allow the production of additional copies of the particle whose position is being measured.
If indeed our Quantum-Gravity scenario hosts a Planck-scale modification of the dispersion relation of the form (1) then clearly the relation between $`\delta p_\gamma `$ and $`\delta E_\gamma `$ should be re-written as follows
$$\delta p_\gamma \left(1+\alpha _1L_pE+3\left(\frac{\alpha _2}{2}\frac{\alpha _1^2}{8}\right)L_p^2E^2\right)\delta E_\gamma $$
(4)
which then leads to the requirement
$$M\frac{1}{\delta x}\alpha _1\frac{L_p}{(\delta x)^2}+\left(\frac{11}{8}\alpha _1^2\frac{3}{2}\alpha _2\right)\frac{L_p^2}{(\delta x)^3}+O\left(\frac{L_p^3}{(\delta x)^4}\right).$$
(5)
These results strictly apply to the measurement of the position of a particle at rest, but they can be straightforwardly generalized (simply using a boost) to the case of the measurement of the position of a particle of energy $`E`$. For the standard case this leads to the $`E1/\delta x`$ relation while in presence of an MDR one easily finds
$$E\frac{1}{\delta x}\alpha _1\frac{L_p}{(\delta x)^2}+\left(\frac{11}{8}\alpha _1^2\frac{3}{2}\alpha _2\right)\frac{L_p^2}{(\delta x)^3}+O\left(\frac{L_p^3}{(\delta x)^4}\right).$$
(6)
### 2.3 GUP and a Planck-scale particle-localization limit
While the connection between a MDR and a Planck-scale particle-localization limit is somewhat less obvious (and in fact we found no mention of it in the literature previous to Ref. ), it is not at all surprising that the GUP would give rise to such a particle-localization limit. In fact, as mentioned, the GUP is primarily viewed as a way to introduce a Planckian limit on the size of the collision region, applicable to high-energy microscopic collision processes, and a limitation on the size of collision regions would naturally be expected to lead to a particle-localization limit. Indeed, as the careful reader can easily verify, from the GUP one obtains (following again straightforwardly the familiar line of analysis discussed in Ref. ) a modification of the relation $`E1/\delta x`$. The modification is of the type $`E1/\delta x+\mathrm{\Delta }`$, with $`\mathrm{\Delta }`$ of order $`\alpha L_p^2/\delta x^3`$, and originates from the fact that according to the GUP, (2), one obtains $`\delta p_\gamma 1/\delta x+\lambda _s^2/\delta x^3`$ (instead of the original $`\delta p_\gamma 1/\delta x`$). Using the standard special-relativistic dispersion relation for a photon $`p_\gamma =E_\gamma `$ the condition on the momentum uncertainty translates in a condition on the energy uncertainty $`\delta E_\gamma \frac{1}{\delta x}\left(1+\alpha \frac{L_p^2}{\delta x^2}\right)`$, and ultimately this leads to
$$E\frac{1}{\delta x}+\alpha \frac{L_p^2}{(\delta x)^3}+O\left(\frac{L_p^3}{(\delta x)^4}\right).$$
(7)
## 3 MDR and black hole entropy
In this section we revisit the argument already proposed in Ref. , suggesting that a Planck-scale modification of the particle-localization limit, of the type (6) or (7), can be used to motivate corrections to the $`S=A/(4L_p^2)`$ area-entropy relation for black holes. We focus here on the case of a MDR, but since the key ingredient is the Planck-scale particle-localization limit, one should expect that, as we confirm explicitly in the next section, the same line of analysis is applicable also to the case in which one takes as starting point a GUP. Since the literature on MDRs is composed both of papers using arguments that are only sufficient to specify the first terms in a series expansion of the MDR at energies below the Planck scale, and some analyses proposing a complete all-order formula for the MDR, we find appropriate to consider these possibilities separately. The power-series analysis will already show that the implications of a MDR can be significant. But the power-series analysis can only be reliably used at energies safely below the Planck scale. In considering some examples of all-order MDR proposals we will also develop some intuition for the type of implications that a MDR could have for black-hole physics at Planckian energy scales.
Our argument connecting a MDR (a particle-localization limit) and some modifications of the area-entropy relation for black holes is formulated in a scheme of analysis first introduced by Bekenstein , which is actually one of the classic arguments for the description of the entropy-area relation. In order to render our presentation self-contained we open this section by describing this classic Bekenstein argument, but we just sketch out the Bekenstein derivation since we expect most readers to be already familiar with it.
### 3.1 The original Bekenstein argument (with unmodified dispersion relation and unmodified uncertainty principle)
The argument presented by Bekenstein in Ref. uses very simple ingredients to suggest that the entropy of a black hole should be proportional to its (horizon-surface) area. The starting point is the general-relativity result establishing that the minimum increase of area when the black hole absorbs a classical particle of energy $`E`$ and size $`s`$ is $`\mathrm{\Delta }A8\pi L_p^2Es`$ (in โnatural unitsโ with $`\mathrm{}=c=1`$). In order to describe the absorption of a quantum particle one must describe the size of the particle in terms of the uncertainty in its position , $`s\delta x`$, and take into account a โcalibration factorโ $`(\mathrm{ln}2)/2\pi `$ that connects the $`\mathrm{\Delta }A8\pi L_p^2Es`$ classical-particle result with the quantum-particle estimate $`\mathrm{\Delta }A4(\mathrm{ln}2)L_p^2E\delta x`$. Bekenstein then enforces the requirement that a particle with position uncertainty $`\delta x`$ should at least have energy $`E1/\delta x`$, which leads to $`\mathrm{\Delta }A4(\mathrm{ln}2)L_p^2`$, and assumes that the entropy depends only on the area of the black hole. Also using the fact that the minimum increase of entropy should be $`\mathrm{ln}2`$, independently of the value of the area, one then concludes that
$$\frac{dS}{dA}\frac{min(\mathrm{\Delta }S)}{min(\mathrm{\Delta }A)}\frac{\mathrm{ln}2}{4(\mathrm{ln}2)L_p^2}.$$
(8)
From this it follows that (up to an irrelevant constant contribution to entropy):
$$S\frac{A}{4L_p^2}.$$
(9)
### 3.2 MDR and black hole entropy in leading order
The Bekenstein argument implicitly assumes (through the $`E1/\delta x`$ relation) that the energy-momentum dispersion relation and the position-momentum uncertainty principle take the standard form. Let us now reformulate the argument, still assuming a standard form for the position-momentum uncertainty principle, but introducing a MDR of the type (1). As in the original Bekenstein argument , we take as starting point the general-relativity result which establishes that the area of a black hole changes according to $`\mathrm{\Delta }A8\pi L_p^2Es`$ when a classical particle of energy $`E`$ and size $`s`$ is absorbed. And again we describe the size of the particle in terms of the uncertainty in its position as done in the previous subsection, obtaining $`\mathrm{\Delta }A4(\mathrm{ln}2)L_p^2E\delta x`$. Whereas in the original Bekenstein argument one then enforces the relation $`E1/\delta x`$ (and this leads to $`\mathrm{\Delta }A4(\mathrm{ln}2)L_p^2`$), we must take into account the MDR-induced Planck-length modification in (6), obtaining
$`\mathrm{\Delta }A`$ $``$ $`4(\mathrm{ln}2)\left[L_p^2{\displaystyle \frac{\alpha _1L_p^3}{\delta x}}{\displaystyle \frac{\left(\frac{3}{2}\alpha _2\frac{11}{8}\alpha _1^2\right)L_p^4}{(\delta x)^2}}\right]`$
$``$ $`4(\mathrm{ln}2)\left[L_p^2{\displaystyle \frac{\alpha _1L_p^3}{R_S}}{\displaystyle \frac{\left(\frac{3}{2}\alpha _2\frac{11}{8}\alpha _1^2\right)L_p^4}{(R_S)^2}}\right]`$
$``$ $`4(\mathrm{ln}2)\left[L_p^2{\displaystyle \frac{\alpha _12\sqrt{\pi }L_p^3}{\sqrt{A}}}{\displaystyle \frac{\left(\frac{3}{2}\alpha _2\frac{11}{8}\alpha _1^2\right)4\pi L_p^4}{A}}\right],`$
where we also used the fact that in falling in the black hole the particle acquires position uncertainty $`\delta xR_S`$, where $`R_S`$ is the Schwarzschild radius (and of course $`A=4\pi R_S^2`$). From (3.2) we derive an area-entropy relation assuming that the entropy of the black hole depends only on its area and that the minimum increase of entropy should be, independently of the value of the area, $`\mathrm{ln}2`$:
$$\frac{dS}{dA}\frac{min(\mathrm{\Delta }S)}{min(\mathrm{\Delta }A)}\frac{\mathrm{ln}2}{4(\mathrm{ln}2)L_p^2\left[1\frac{\alpha _12\sqrt{\pi }L_p}{\sqrt{A}}\frac{\left(\frac{3}{2}\alpha _2\frac{11}{8}\alpha _1^2\right)4\pi L_p^2}{A}\right]}\left(\frac{1}{4L_p^2}+\frac{\alpha _1\sqrt{\pi }}{2L_p\sqrt{A}}+\frac{\left(\frac{3}{2}\alpha _2\frac{11}{8}\alpha _1^2\right)\pi }{A}\right),$$
(11)
which gives (up to an irrelevant constant contribution to entropy)
$$S\frac{A}{4L_p^2}+\alpha _1\sqrt{\pi }\frac{\sqrt{A}}{L_p}+\left(\frac{3}{2}\alpha _2\frac{11}{8}\alpha _1^2\right)\pi \mathrm{ln}\frac{A}{L_p^2}.$$
(12)
This result of course reproduces the famous linear formula if all coefficients $`\alpha _i`$ vanish. If the cubic term $`\alpha _1E^3`$ is present in the energy-momentum dispersion relation then the leading correction goes like $`\sqrt{A}`$, whereas if the first nonzero coefficient in the dispersion relation expansion is $`\alpha _2`$ the leading correction term goes like $`\mathrm{log}A`$. Our โimproved Bekenstein argumentโ therefore provides a possible link between the form of the MDR (and of the GUP, as we stress later) and the all-order form of the entropy-area relation for black holes. For example, if within a given quantum-gravity approach one can find a general argument suggesting that there are no $`\sqrt{A}`$ terms in the entropy-area relation, then one can use our improved Bekenstein argument to deduce that within that given quantum-gravity approach one should not find terms of the type $`\alpha _1E^3`$ in the energy-momentum dispersion relation.
Over the last few years both in String Theory and in Loop Quantum Gravity some techniques for the direct analysis of the entropy of black holes, using their quantum properties, have been developed, and these techniques are now able to go even beyond the entropy-area-proportionality contribution: they establish that the leading correction should be of log-area type, so that one expects (for $`AL_p^2`$) an entropy-area relation for black holes of the type
$$S=\frac{A}{4L_p^2}+\rho \mathrm{ln}\frac{A}{L_p^2}+O\left(\frac{L_p^2}{A}\right).$$
(13)
where $`\rho `$ is a coefficient which might take different value in String Theory and in Loop Quantum Gravity. The status of the energy-momentum dispersion relation within these theories is not completely settled (it is much debated particularly in the Loop-Quantum-Gravity literature), but on the basis of our improved Bekenstein argument we can conclude that both String Theory and Loop Quantum Gravity cannot allow terms of the type $`\alpha _1E^3`$ in the energy-momentum dispersion relation. And, if their log-area corrections to the entropy-area relation are to be trusted, we expect that both String Theory and Loop Quantum Gravity should either predict a $`\alpha _2E^4`$ correction to the dispersion relation or (see later) they should host a corresponding modification of the position-momentum uncertainty principle.
We can also use (12) to obtain, using the first law of black hole thermodynamics $`dS=\frac{dM}{T}`$, a Planck-scale-corrected relation between black-hole temperature and mass:
$$T_{BH}^{MDR}\frac{E_p^2}{8\pi M}\left(1\alpha _1\frac{E_p}{2\sqrt{2}M}\left(\frac{15}{32}\alpha _1^2\frac{3}{8}\alpha _2\right)\frac{E_p^2}{M^2}\right),$$
(14)
where we also used the familiar relation between black hole area and mass $`A=16\pi M^2`$.
### 3.3 Some all-order results for MDR modifications of black-hole entropy
In the previous subsection we were establishing a possible relation between MDR and log corrections to the entropy-area relation. Since the log-area term is a leading-order term it was appropriate to work within a power-series expansion of the MDR. Moreover, the mentioned results from quantum-gravity research (primarily from Loop Quantum Gravity and approaches based on noncommutative geometry) that provide motivation for a Planck-scale modification of the dispersion relation in most cases are obtained within analyses that only have access to the first terms in a power-series expansion of the dispersion relation. Still for some aspects of our analysis it will be useful to contemplate some illustrative examples of all-order dispersion relations, especially when we try to figure out what could be some examples of implications of a MDR for the behaviour of black holes of Planck-length size.
The careful reader can easily verify that once a given energy-momentum dispersion relation $`E=f_{disp}(p)`$ is adopted the steps of the calculation reported in the preceding subsection can be followed rather straightforwardly, obtaining
$$\frac{dS}{dA}\frac{min(\mathrm{\Delta }S)}{min(\mathrm{\Delta }A)}\frac{1}{2L_p^2}\sqrt{\frac{\pi }{A}}\frac{1}{f_{disp}\left(\sqrt{\frac{4\pi }{A}}\right)}$$
(15)
and
$$T_{BH}\frac{1}{4\pi }f_{disp}\left(\frac{E_p^2}{2M}\right).$$
(16)
As illustrative examples of โall-order MDRsโ we consider the following three cases:
$`\mathrm{cosh}(E/E_p)\mathrm{cosh}(m/E_p){\displaystyle \frac{p^2}{2E_p^2}}e^{E/E_p}`$ $`=`$ $`0,`$ (17)
$`{\displaystyle \frac{E^2}{(1E/E_p)^2}}{\displaystyle \frac{p^2}{(1E/E_p)^2}}m^2`$ $`=`$ $`0,`$ (18)
$`\mathrm{cosh}(\sqrt{2}E/E_p)\mathrm{cosh}(\sqrt{2}m/E_p){\displaystyle \frac{p^2}{E_p^2}}\mathrm{cosh}(\sqrt{2}E/E_p)`$ $`=`$ $`0,`$ (19)
(17) has already been considered in the previous literature , particularly as a possible description of particle propagation in $`\kappa `$-Minkowski noncommutative spacetime. It provides an example in which the coefficient of the linear-in-$`L_p`$ term is nonvanishing: $`\alpha _1=1/2`$. And it is noteworthy that according to (17) there is a maximum momentum for fundamental particles: from (17) it follows that for $`E\mathrm{}`$ one has $`pE_p`$.
The case (19) has not been previously considered in the literature. It provides for our purposes a valuable illustrative example since, as in the case of (17), it would lead to a maximum momentum ($`pE_p`$ for $`E\mathrm{}`$) but, contrary to the case of (17), it corresponds to $`\alpha _1=0`$ (whereas $`\alpha _2=5/18`$). This is therefore an example with the maximum-momentum feature and such that one would expect the leading corrections to the entropy-area relation to be logarithmic.
The case of (18) has already been considered in the literature for other reasons , and it provides us an opportunity to illustrate some consequences of a scenario in which both $`\alpha _1`$ and $`\alpha _2`$ vanish, but still there are some Planck-scale modifications of the energy-momentum dispersion relation. And it is noteworthy that (18) can be implemented in such a way that the Planck scale provides the maximum value of both momentum and energy.
For the cases with dispersion relations (17) or (19), since $`E\mathrm{}`$ for $`pE_p`$, the formulas derived above would lead to the conclusion that the black hole temperature diverges at some finite (nonzero!) value of the black-hole mass $`M_{min}=E_p/2`$. We would then assume that this $`M_{min}`$ is the minimum allowed mass for a black hole, and that the standard description of the evaporation process should not be applicable beyond this small value of mass.
In cases in which one introduces both a maximum momentum and a maximum energy while keeping the form of the dispersion relation largely unaffected<sup>3</sup><sup>3</sup>3Whenever the mass $`m`$ can be ignored (i.e. for massless particles and high-energy particles with finite mass) the dispersion relation (18) is indistinguishable from the standard special-relativistic one., as done in some applications of (18), one would expect (since the energy has a maximum Planckian value, $`E_{Max}=E_P`$) that the temperature should be bounded to be lower than the Planck scale, $`T_{Max}E_p`$, and that the minimum allowed value of black-hole mass should be also Plankian, since it should be the value of mass such that temperature reaches is maximum allowed value.
## 4 GUP (with and without MDR) and black hole entropy
In the previous section we focused on scenarios in which the energy-momentum dispersion relation is modified but the position-momentum uncertainty principle preserves the Heisenberg form. But clearly the key ingredient of our analysis is the presence of a correction term $`\mathrm{\Delta }`$ in the particle-localization-limit relation $`E1/\delta x+\mathrm{\Delta }`$. As stressed in Section 2, both a MDR and a GUP can introduce such a correction term in the particle-localization limit, and therefore, as we want to discuss explicitly in this Section, also in presence of a GUP one should expect corrections to the entropy-area black-hole formula and to the formula that relates the mass and the temperature of a black hole.
We start the section by considering scenarios in which the position-momentum uncertainty principle is Planck-scale modified, while the energy-momentum dispersion relation preserves its special-relativistic form. Then in Subsection 4.2 we comment on the more general case, in which one might be dealing with both a MDR and a GUP.
### 4.1 GUP and black hole entropy
Let us start by noting here again for convenience the particle-localization limit that one obtains assuming a GUP of the form (2) and a standard (special-relativistic) energy-momentum dispersion relation:
$$E\frac{1}{\delta x}+\alpha \frac{L_p^2}{(\delta x)^3}+O\left(\frac{L_p^3}{(\delta x)^4}\right).$$
(20)
Following the same strategy of analysis adopted in the previous section, one finds that the Bekenstein argument, when taking into account this localization limit (20), leads to the conclusion that the maximum increase of black-hole area upon absorption of a particle of energy $`E`$ is given by
$`\mathrm{\Delta }A4(\mathrm{ln}2)\left[L_p^2+{\displaystyle \frac{\alpha L_p^4}{(\delta x)^2}}\right]4(\mathrm{ln}2)\left[L_p^2+{\displaystyle \frac{\alpha L_p^4}{(R_S)^2}}\right]4(\mathrm{ln}2)\left[L_p^2+{\displaystyle \frac{\alpha 4\pi L_p^4}{A}}\right].`$
From this it follows that the entropy-area relation should take the form
$$S\frac{A}{4L_p^2}\alpha \pi \mathrm{ln}\frac{A}{L_p^2},$$
(21)
and the formula relating the temperature and the mass of the black hole should take the form
$$T_{BH}^{GUP}\frac{E_p^2}{8\pi M}\left(1+\alpha \frac{E_p^2}{8M^2}\right).$$
(22)
### 4.2 Combining MDR and GUP in the analysis of black-hole entropy
We have argued that both a MDR and a GUP are possible features of a quantum-gravity theory that would affect black-hole termodynamics. Actually, as the careful reader must have noticed, the line of analysis we are advocating is composed of two steps. First we notice that the โparticle-localization limitโ in its standard form, $`E1/\delta x`$, is derived on the basis of two key assumptions, the validity of the Heisenberg position-momentum uncertainty principle and the validity of the special-relativistic energy-momentum dispersion relation, and that by modifying the uncertainty principle and/or the dispersion relation one gets a modified particle-localization limit of the type $`E1/\delta x+\mathrm{\Delta }_{\delta x,L_p}`$. Then we observe that a key assumption of the Bekenstein argument for the derivation of black-hole entropy is the validity of the standard particle-localization limit $`E1/\delta x`$. With a MDR and/or a GUP one gets a modified particle-localization limit, which in turn leads to a modification of the black-hole area-entropy relationship.
It is worth mentioning that the modifications induced by a MDR and a GUP may (at least in part) cancel out at the level of the area-entropy equation. In order to stress the importance of this possibility let us consider the information presently available on the Loop Quantum Gravity approach: (i) several Loop-Quantum-Gravity studies have argued in favour of a MDR with nonvanishing $`\alpha _1`$ (leading Planck-scale correction to the dispersion relation that goes linearly with $`L_p`$), (ii) there is no mention of a GUP in the Loop-Quantum-Gravity literature, (iii) several Loop-Quantum-Gravity studies have argued in favour of an entropy-area relationship in which the leading correction, beyond the linear term, is of log-area type. According to the perspective on the derivation of black-hole entropy that we are advocating one would find these three ingredients to be logically incompatible: if the MDR has nonvanishing $`\alpha _1`$ and the position-momentum uncertainty principle is not Planck-scale modified then in the entropy-area relationship the leading correction, beyond the linear term, should have $`\sqrt{area}`$ dependence. Does this mean that Loop Quantum Gravity is a logically inconsistent framework? Of course, it does not. It simply means that some of the relevant preliminary results must be further investigated. It may well be that, as the loop-quantum-gravity approach is understood more deeply, it turns out that the $`\alpha _1`$ coefficient in the MDR vanishes. Or else we might discover that in Loop Quantum Gravity the $`\alpha _1`$ coefficient in the MDR takes a nonzero value, but there is a corresponding linear-in-$`L_p`$ term in the GUP with just the right coefficient to give an overall vanishing coefficient to the $`\sqrt{area}`$ term in the entropy-area relation.
Our perspective on the derivation of black-hole entropy provides a logical link between different aspects of a quantum-gravity theory and may be used most fruitfully when, as in the case of Loop Quantum Gravity, the formalism is very rich and some of the results obtained within that formalism are of preliminary nature. Even before being able to derive more robust results we may uncover that the presently-available preliminary results are not providing us with a logically-consistent picture, and this in turn will give us additional motivation for investigating more carefully those preliminary results.
It is also worth mentioning that on the string-theory side our perspective on the derivation of black-hole entropy provides no evidence of a logical inconsistency among the results so far obtained in that framework. The string-theory literature indicates that the entropy-area relationship should involve a leading correction, beyond the linear term, of log-area type, and provides strong evidence of a GUP of the type (2), while the results so far obtained do not indicate the need to modify the dispersion relation in string theory. These three ingredients provide a logically-consistent scenario within our perspective on the derivation of black-hole entropy. As shown above, with a GUP of the type (2) and with an unmodified (still special-relativistic) dispersion relation one is indeed led to an entropy-area relationship in which the leading correction, beyond the linear term, is of log-area type.
## 5 Implications for the Bekenstein entropy bound and Generalized Second Law
It is natural at this point, after having shown that a MDR and a GUP can affect the black-hole entropy-area and mass-temperature relationships, to wonder whether other aspects of black-hole thermodynamics are also affected, and whether the overall picture preserves the elegance/appeal of the original scheme, based on standard uncertainty principle and dispersion relation. In this section we investigate the validity of the Generalized Second Law (GSL) of thermodynamics and the implications for the Bekenstein entropy bound. In order to work within a definite scenario we assume here a MDR (while we implicitly assume that the uncertainty principle takes its standard form).
The GSL asserts that the second law of thermodynamics is still valid in presence of collapsed matter. Given the entropy of the black hole, as described by the area-entropy relation, the GSL requires that the total entropy of a system composed of a black hole and ordinary matter never decreases. This means that the following inequality holds for all physical processes
$$S_{BH}+S_{mat}0.$$
(23)
It was observed that in principle (using the so-called โGeroch processโ) one could violate the GSL if objects of fixed size $`R`$ and energy $`E`$ could have arbitrarily large entropy $`S`$. This led Bekenstein to propose a โentropy boundโ
$$S_{mat}2\pi ER$$
(24)
for an arbitrary system of energy $`E`$ and effective radius $`R`$. The fact that the GSL implies the Bekenstein bound and vice versa has long been debated and is still actively debated. However the Bekenstein bound turns out to hold for a variety of systems in flat Minkowski space and can be derived as weak-gravity limit of the popular โGeneralized Covariant Entropy Boundโ .
A remarkable feature of the Bekenstein bound is that, in spite of being motivated by considerations rooted in the gravitational realm, it does not involve the Planck scale (or equivalently Newtonโs constant). The absence of the Planck scale is less puzzling in light of the observation that the bound can be derived even without advocating gravity in any way: it is sufficient to analyze some implications of the particle-localization limit $`E\frac{1}{\delta x}`$. This alternative derivation requires considering a matter system with energy $`E`$, in which self-gravitation effects can be neglected, that occupies a region in flat spacetime with radius $`R`$ smaller than the gravitational radius $`R_G2L_p^2E`$. The standard particle-localization limit, when generalized to this type of systems, sets a minimum value for the energy of a quantum in a region of spatial radius $`R`$
$$ฯต(R)\frac{1}{R}.$$
(25)
The maximum number of quanta that we can have in the region is then given by
$$N_{max}\frac{E}{ฯต(R)}=ER.$$
(26)
If we consider the simple case of a system for which the maximal number of microstates for $`N`$ particles is given by $`\mathrm{\Omega }(N)=2^N`$ then the entropy of the system $`S=log\mathrm{\Omega }(N)`$ is bounded by the inequality
$$S_{mat}\mathrm{log}2ER,$$
(27)
which is indeed consistent with the Bekenstein bound (up to another โcalibration factorโ $`\eta =\frac{2\pi }{log2}`$).
We briefly reviewed this derivation of the Bekenstein bound especially in order to stress the role played by the particle-localization limit $`E\frac{1}{\delta x}`$. It is then obvious that the modifications of the particle-localization limit induced by a MDR (and/or a GUP) would affect the Bekenstein bound. As shown earlier, within our parametrization of the MDR<sup>4</sup><sup>4</sup>4For an analogous modification of the Bekenstein bound coming from the GUP see Ref. ., one obtains a particle localization limit of the form
$$ฯต(R)\frac{1}{R}\left(1\alpha _1\frac{L_p}{R}\left(\frac{3}{2}\alpha _2\frac{11}{8}\alpha _1^2\right)\frac{L_p^2}{R^2}+O\left(\frac{L_p^3}{R^3}\right)\right)$$
(28)
which gives
$$S_{mat}2\pi ER\left(1+\alpha _1\frac{L_p}{R}+\left(\frac{3}{2}\alpha _2\frac{11}{8}\alpha _1^2\right)\frac{L_p^2}{R^2}+O\left(\frac{L_p^3}{R^3}\right)\right).$$
(29)
This MDR-modified Bekenstein bound fits very naturally with our corresponding formula, (12), for the entropy-area relation; in fact, the two results combine to provide us with a picture which is still consistent with the GSL. According to (29) when a matter system of energy $`E`$ falls into the black hole, this corresponds to a negative change of entropy which has absolute value not greater than
$$max(|\mathrm{\Delta }S_{mat}|)2\pi ER\left(1+\alpha _1\frac{L_p}{R}+\left(\frac{3}{2}\alpha _2\frac{11}{8}\alpha _1^2\right)\frac{L_p^2}{R^2}+O\left(\frac{L_p^3}{R^3}\right)\right)$$
(30)
and correspondingly, according to (12), the black hole entropy increases at least by
$$min(\mathrm{\Delta }S_{BH})2\pi ER\left(1+\alpha _1\frac{L_p}{R}+\left(\frac{3}{2}\alpha _2\frac{11}{8}\alpha _1^2\right)\frac{L_p^2}{R^2}+O\left(\frac{L_p^3}{R^3}\right)\right)$$
(31)
Thus the MDR-induced corrections to $`S_{BH}`$ and $`S_{mat}`$ cancel exactly at the level of the inequality relevant for the GSL. The GSL stills holds, even in presence of a modified particle-localization limit.
## 6 Corrections to black-body radiation spectrum
In preparation for some observations on black-hole evaporation, to which we devote Section 7, we now want to investigate the implications of a MDR and/or a GUP for the black-body radiation spectrum.
### 6.1 MDR and black-body spectrum in leading order
Let us start by considering photons in a cubical box with edges of length $`L`$ (and volume $`V=L^3`$). The wavelengths of the photons are subject to the boundary condition $`\frac{1}{\lambda }=\frac{n}{2L}`$, where $`n`$ is a positive integer. This condition implies, assuming that the de Broglie relation is left unchanged, that the photons have (space-)momenta that take values $`p=\frac{n}{2L}`$. Thus momentum space is divided into cells of volume $`V_p=\left(\frac{1}{2L}\right)^3=\frac{1}{8V}`$. From this it follows that the number of modes with momentum in the interval $`[p,p+dp]`$ is given by
$$g(p)dp=8\pi Vp^2dp.$$
(32)
Assuming a MDR of the type parametrized in (1) one then finds that ($`m=0`$ for photons)
$$pE\left(1+\frac{\alpha _1}{2}L_pE+\left(\frac{\alpha _2}{2}\frac{\alpha _1^2}{8}\right)L_p^2E^2\right)$$
(33)
and
$$dp\left(1+\alpha _1L_pE+\left(\frac{3}{2}\alpha _2\frac{3}{8}\alpha _1^2\right)L_p^2E^2\right)dE$$
(34)
Using this in (32) one obtains
$$g(E)dE=8\pi V\left(1+2\alpha _1L_pE+5\left(\frac{1}{2}\alpha _2+\frac{1}{8}\alpha _1^2\right)L_p^2E^2\right)E^2dE$$
(35)
which in terms of the frequency $`\nu `$ takes the form
$$g(\nu )d\nu =8\pi V\left(1+2\alpha _1L_p\nu +5\left(\frac{1}{2}\alpha _2+\frac{1}{8}\alpha _1^2\right)L_p^2\nu ^2\right)\nu ^2d\nu .$$
(36)
In order to obtain the MDR-modified energy density of a black body at temperature $`T`$ we must now use (36) and rely on the statistical arguments which show that in a system of bosons at temperature $`T`$ the average energy per oscillator is given by
$$\overline{E}=\frac{\nu }{e^{\frac{\nu }{T}}1}.$$
(37)
Thus the energy density at a given temperature $`T`$, for the frequency interval $`[\nu ,\nu +d\nu ]`$, is
$$u_\nu (T)d\nu =8\pi \left(1+2\alpha _1L_p\nu +5\left(\frac{1}{2}\alpha _2+\frac{1}{8}\alpha _1^2\right)L_p^2E^2\right)\frac{\nu ^3d\nu }{e^{\frac{\nu }{T}}1}.$$
(38)
and integrating this formula we get the MDR-modified energy density of a black body at temperature $`T`$
$$u(T)=\frac{8\pi ^5}{15}T^4+384\pi \zeta (5)\alpha _1L_pT^5+5\left(\frac{1}{2}\alpha _2+\frac{1}{8}\alpha _1^2\right)\frac{160\pi ^7}{63}L_p^2T^6$$
(39)
The MDR introduces corrections of the type $`T^{4+n}/E_P^n`$ to the Stefan-Boltzmann law. Moreover, the maximum value of the integrand in (38), as a function of $`\nu `$, is clearly also shifted: the MDR also introduces a modification of Wienโs law. Of course, using the low-energy expansion (1) of the dispersion relation we only get a reliable picture at temperatures safely below the Planck scale, but the presence of correction terms of the type $`T^{4+n}/E_P^n`$ clearly suggests that the MDR-modified description leads to departures from the Stefan-Boltzmann law that can become very significant as the temperature approaches the Planck scale. We intend to show this explicitly by considering an example of all-order MDR formula.
### 6.2 Some all-order results for MDR modifications of black-body spectrum
Let us therefore derive once again the modified Stefan-Boltzmann law, now assuming, as illustrative example of an all-order MDR formula, the validity of the dispersion relation (19). Clearly the number of modes in momentum space is still given by
$$g(p)dp=8\pi Vp^2dp,$$
(40)
but now
$$p^2=E_p^2\left(1\frac{1}{\mathrm{cosh}(\sqrt{2}E/E_p)}\right)$$
(41)
and this implies that the number of modes for given energy is given by
$$g(E)dE=16\pi VE_p^2\mathrm{sinh}^2\left(\frac{E/E_p}{\sqrt{2}}\right)\mathrm{cosh}\left(\frac{E/E_p}{\sqrt{2}}\right)\frac{1}{\mathrm{cosh}^{5/2}\left(\sqrt{2}E/E_p\right)}dE$$
(42)
i.e. the number of modes for given frequency is
$$g(\nu )d\nu =16\pi VE_p^2\mathrm{sinh}^2\left(\frac{\nu /E_p}{\sqrt{2}}\right)\mathrm{cosh}\left(\frac{\nu /E_p}{\sqrt{2}}\right)\frac{1}{\mathrm{cosh}^{5/2}\left(\sqrt{2}\nu /E_p\right)}d\nu .$$
(43)
Then the modified Stefan-Boltzmann law is given, in integral form, by
$$u(T)=\frac{1}{V}_0^{\mathrm{}}\frac{g(\nu )}{e^{\frac{\nu }{T}}1}\nu ๐\nu ,$$
(44)
where $`g(\nu )`$ is the one of (43).
It is useful to consider some limiting forms of the integration in (44). Clearly, since (19) is consistent with (1) for $`\alpha _1=0`$ and $`\alpha _2=5/18`$, in the limit $`T/E_p1`$ the integration (44) gives a result that reproduces (39) for $`\alpha _1=0`$ and $`\alpha _2=5/18`$. But, now that we are dealing with an all-order formula, besides considering the case $`T/E_p1`$ we can also investigate the opposite limit $`T/E_p1`$, finding
$$u(T)=16\pi E_p^4\left\{\frac{T}{E_p}C_1\frac{1}{2}C_2\frac{E_p}{T}C_3+O(E_p^2/T^2)\right\}$$
(45)
where
$`C_1`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}\mathrm{sinh}^2(x/\sqrt{2}){\displaystyle \frac{\mathrm{cosh}(x/\sqrt{2})}{\mathrm{cosh}^{5/2}(\sqrt{2}x)}}๐x={\displaystyle \frac{1}{6}},`$ (46)
$`C_2`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}x\mathrm{sinh}^2(x/\sqrt{2}){\displaystyle \frac{\mathrm{cosh}(x/\sqrt{2})}{\mathrm{cosh}^{5/2}(\sqrt{2}x)}}๐x0.22,`$ (47)
$`C_3`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}x^2\mathrm{sinh}^2(x/\sqrt{2}){\displaystyle \frac{\mathrm{cosh}(x/\sqrt{2})}{\mathrm{cosh}^{5/2}(\sqrt{2}x)}}๐x0.41,`$ (48)
This means that the MDR (19) leads to a modification of the Stefan-Boltzmann law which at the Planck scale is very significant: for $`TE_p`$ one finds that $`u`$ depends linearly on $`T`$, rather than with the fourth power.
It is of particular interest to establish what is the relationship between the โcharacteristic frequencyโ (and characteristic wavelength) of the black-body spectrum and temperature. In the standard description of a black body the characteristic frequency grows linearly with the temperature. In order to verify whether this is still the case in our MDR-modified scenario we can take the derivative of $`u_\nu (T)`$ with respect to $`\nu `$, so that we can identify the value of frequency for which the energy density (and the radiated flux) reaches a maximum. This leads to the following equation that must be satisfied by the characteristic frequency $`\overline{\nu }`$:
$$\left(e^{\frac{\overline{\nu }}{T}}1\right)\left(g(\overline{\nu })+g^{}(\overline{\nu })\overline{\nu }\right)\frac{e^{\frac{\overline{\nu }}{T}}}{T}g(\overline{\nu })\overline{\nu }=0.$$
(49)
For $`TE_p`$ of course this reproduces the type of small modification of Wienโs law, which we already noticed in the previous section. The fact that we are now considering a scenario with a given all-order MDR formula allows us to examine the dependence of the characteristic frequency on temperature even when the temperature reaches and eventually exceeds the Planck scale. And we find that the for $`T/E_p1`$ the characteristic frequency becomes essentially independent of temperature. No matter how high the temperature goes the characteristic frequency never exceeds the following finite value:
$$\overline{\nu }E_p\frac{\mathrm{cosh}^1[(1+\sqrt{41})/4]}{\sqrt{2}}0.87E_p$$
(50)
So basically at low temperatures any increase of temperature causes a corresponding increase in characteristic frequency of the black-body spectrum, but gradually a saturation mechanism takes over and even in the infinite-temperature limit the characteristic frequency is still finite, and given by the Planck scale (up to a coefficient of order 1). This occurs with the dispersion relation (19), i.e. in a scenario with a minimum value of wavelength but no maximum value of frequency. An analogous result for the case of the dispersion relation (18), which leads to both a minimum value of wavelength and a maximum value of frequency, would have not been surprising: if the framework introduces from the beginning a maximum Planckian value of frequency, then of course also the characteristic frequency of black-body radiation would be โsubPlanckianโ. But in analyzing the case of (19) we found that the presence of a minimum wavelength at the fundamental level is sufficient for the emergence of a maximum Planckian value of the characteristic frequency of black-body radiation, as shown explicitly by Eq. (50).
### 6.3 Black-body spectrum with GUP
In the previous two subsections the key point was that a MDR leads to a modified formula for the density of modes in a given (infinitesimal) frequency interval, $`g(\nu )d\nu `$. If instead we now assume that the dispersion relation takes its standard special-relativistic form, but there is a GUP, it is not a priori obvious that the black-body spectrum is affected. One does indeed obtain a modified black-body spectrum if it is assumed that the GUP should also be reflected in a corresponding modification of the de Broglie relation,
$$\lambda \frac{1}{p}\left(1+\alpha L_p^2p^2\right)$$
(51)
and
$$E\nu \left(1+\alpha L_p^2\nu ^2\right).$$
(52)
For oscillators in a box the number of modes in an infinitesimal frequency interval would still be described by the standard formula
$$g(\nu )d\nu =8\pi V\nu ^2d\nu ,$$
(53)
but, as a result of (52), the average energy per oscillator would be given by
$$\overline{E}=\frac{\nu }{e^{\frac{\nu }{T}}1}\left(1+\alpha L_p^2\nu ^2\left(1\frac{\frac{\nu }{T}}{1e^{\frac{\nu }{T}}}\right)\right).$$
(54)
Combining (52) and (54) one finds
$$u_\nu (T)d\nu =8\pi \left(1+\alpha L_p^2\nu ^2\left(1\frac{\frac{\nu }{T}}{1e^{\frac{\nu }{T}}}\right)\right)\frac{\nu ^3d\nu }{e^{\frac{\nu }{T}}1}.$$
(55)
and the modified Stefan-Boltzmann law takes the form
$$u(T)=\frac{8\pi ^5}{15}T^4+\frac{8\pi ^6}{9}\alpha L_p^2T^6.$$
(56)
The $`L_p^2T^6`$ correction term is just one of the $`L_p^nT^{4+n}`$ correction terms on which we already commented in the context of the MDR modifications of black-body radiation.
## 7 Black hole evaporation
In this section we use some of the results obtained in the previous sections in a description of the black-hole evaporation process. The key ingredients are the relation between the black-hole temperature and mass and the relation between the black-hole temperature and the energy density emitted by the black hole. We neglect possible non-thermal corrections due to back-reaction effects (see, e.g., the recent studies in Ref. and references therein), and we therefore treat the radiation emitted by the black-hole as black-body radiation.
### 7.1 MDR and Black hole evaporation
At temperature $`T`$ the intensity $`I`$ of the radiation emitted by a black hole of area $`A`$ is given by
$$I(T)=Au(T).$$
(57)
Using energy conservation one can write
$$\frac{dM}{dt}=Au,$$
(58)
and assuming a MDR of the type $`E=f_{disp}(p)`$, in light of our result (16), one finds
$$\frac{dM}{dt}=16\pi \frac{M^2}{E_p^4}u\left(\frac{1}{4\pi }f_{disp}\left(\frac{E_p^2}{2M}\right)\right)$$
(59)
When $`ME_p`$ (so that a power-series expansion of $`f_{disp}(E_p^2/2M)`$ is meaningful) this takes the form
$`{\displaystyle \frac{dM}{dt}}=`$ $`=`$ $`k_0{\displaystyle \frac{E_p^8}{M^2}}k_1\alpha _1{\displaystyle \frac{E_p^9}{M^3}}(k_{21}\alpha _1^2+k_{22}\alpha _2){\displaystyle \frac{E_p^{10}}{M^4}}+O(E_p^5/M^5)`$ (60)
where $`k_0=\frac{\pi ^2}{480}`$, $`k_1=k_0\frac{90\zeta (5)\pi ^5}{\pi ^5}`$, $`k_{21}=k_0\frac{502\pi ^575600\zeta (5)}{672\pi ^5}`$ and $`k_{22}=k_0\frac{211}{672\pi ^5}`$
This power-series analysis allows to conclude that a MDR can affect the speed of evaporation of a black hole. For example, in the case of the dispersion relation (17) the evaporation process is retarded with respect to the standard case, whereas in the case of (19) the evaporation process is accelerated.
With a given all-order MDR formula one can obtain of course even more detailed information than available using the power-series expansion. In particular, let us look at the case of the dispersion relation (19) and analyze the stage of the evaporation process when the mass of the black hole is of the order of the Planck scale. For $`ME_P`$ we can approximate the MDR (19) as follows
$$E\frac{E_p}{\sqrt{2}}\mathrm{ln}\left(\frac{2}{1(p/E_p)^2}\right)$$
(61)
and then one finds
$$\frac{dM}{dt}(16\pi )^2M^2\left\{\frac{C_1}{4\pi \sqrt{2}}\mathrm{ln}\left(\frac{2}{1(\frac{E_p}{2M})^2}\right)\frac{1}{2}C_2\right\}.$$
(62)
This shows that, in the case of the MDR (19), the energy flux emitted by the black hole would formally diverge as the black-hole mass approaches $`E_p/2`$. This is mainly a consequence of the fact that the black-hole temperature diverges when $`ME_p/2`$. In the standard description of black-hole evaporation these divergences occur as $`M0`$.
### 7.2 GUP and Black hole evaporation
The observations reported in the previous subsection for the case of a MDR (with unmodified energy-momentum uncertainty relation) can be easily adapted to the complementary situation with a GUP and a standard (unmodified) dispersion relation. One must however assume, as already stressed in Subsection 6.3, that the GUP is reflected in a corresponding modification of the de Broglie relation ($`\lambda (1+\alpha L_p^2p^2)/p`$). In this hypothesis one easily finds that the black hole should lose its mass at a rate given by
$$\frac{dM}{dt}=Au\left(\frac{E_p^2}{2M}\right)=16\pi \frac{8\pi ^5}{15}\left(T\left(\frac{E_p^2}{2M}\right)\right)^4+\frac{8\pi ^6}{9}\alpha L_p^2\left(T\left(\frac{E_p^2}{2M}\right)\right)^6.$$
(63)
Expanding for $`M/E_p1`$ we obtain
$$\frac{dM}{dt}16\pi \frac{E_p^4}{M^2}\left(\stackrel{~}{k}_0+\alpha \stackrel{~}{k}_1\frac{E_p^2}{M^2}\right),$$
(64)
with $`\stackrel{~}{k}_0=\frac{\pi }{7680}`$ and $`\stackrel{~}{k}_1=\frac{1}{294912}+\frac{\pi }{15360}`$.
Clearly the modifications to the black hole evaporation formula obtained in the GUP scenario are qualitatively the same as in the MDR scenario with $`\alpha _1=0`$.
## 8 A possible dependence on the speed law for photons
Throughout our analysis we have implicitly assumed that the law $`v_\gamma =1`$ describing the speed of photons is not affected by the MDR and/or the GUP. The possibility of modifications of the speed law for photons has been however considered rather extensively, particularly in the MDR literature. While several authors have argued that the law $`v_\gamma =1`$ should not be modified even in presence of an MDR (see, e.g., Refs. and references therein), one also finds support in the literature for the proposal (see, e.g., Ref. and references therein) of the law $`v_\gamma =[dE/dp]_{m=0}=[df_{disp}(p)/dp]_{m=0}`$ and the proposal (see, e.g., Ref. and references therein) of the law $`v_\gamma =p/E`$.
For our analysis a key point is that if, instead of $`v_\gamma =1`$, one took $`v_\gamma =[dE/dp]_{m=0}`$ or $`v_\gamma =p/E`$ then the speed of photons would acquire an energy dependence which should be taken into account in some aspects of our derivations. We postpone to future studies this more general analysis, but in order to explore the type of modifications which could be induced by such an energy dependence of the speed of photons we do intend to consider here the description of black-body radiation with the dispersion relation (19), assuming that the speed of photons is governed by either $`v_\gamma =[dE/dp]_{m=0}`$ or $`v_\gamma =p/E`$.
We focus on the emitted โflux densityโ
$$I_\nu =Au_\nu v_\gamma (\nu )$$
(65)
where $`A`$ is the area of the radiating surface and $`u_\nu `$ is the energy density at a given frequency.
Taking $`v_\gamma =p/E`$, from (19) it follows that
$$v_\gamma (\nu )=\frac{p}{E}=\frac{E_p}{E}\sqrt{1\frac{1}{\mathrm{cosh}\left(\frac{\sqrt{2}E}{E_p}\right)}}.$$
(66)
From this it would then follow that the energy flux density is given by
$$I_\nu (T)=4\pi A\sqrt{2}E_p^3\frac{1}{e^{\nu /T}1}\frac{\mathrm{sinh}(\sqrt{2}E/E_p)}{\mathrm{cosh}^3(\sqrt{2}E/E_p)}\left[\mathrm{cosh}(\sqrt{2}E/E_p)1\right].$$
(67)
This suggests that, although there are some small quantitative differences, the qualitative features of black-body radiation with the dispersion relation (19) are largely independent of the choice between $`v_\gamma =1`$ and $`v_\gamma =p/E`$. In particular, from (67) with one finds that the typical frequency of the photons contributing to the energy flux saturates at
$$\overline{\nu }0.76E_p,$$
(68)
which is not much different from the typical frequency found for the case $`v_\gamma =1`$. The analysis of the total emitted energy ($`_0^{\mathrm{}}I_\nu (T)๐\nu `$) also leads to rather small differences between the choices $`v_\gamma =1`$ and $`v_\gamma =p/E`$. In particular from (67) one finds
$$I/A=\frac{8}{15}\pi ^5T^4\left\{1+C_1\left(\frac{T}{E_p}\right)^2+C_2\left(\frac{T}{E_p}\right)^4+O\left(\frac{T}{E_p}\right)^6\right\},$$
(69)
in the limit $`T/E_p1`$, and
$$I/A=E_p^4\left\{\stackrel{~}{C}_1\frac{T}{E_p}+\stackrel{~}{C}_2+\stackrel{~}{C}_3\frac{E_p}{T}+O\left(\frac{E_p}{T}\right)^2\right\},$$
(70)
in the limit $`T/E_p1`$, where $`C_1=\frac{100\pi ^2}{21}`$, $`C_2=\frac{164\pi ^4}{5}`$ and $`\stackrel{~}{C}_1=5.57,\stackrel{~}{C}_1=\pi `$ and $`\stackrel{~}{C}_3=0.79`$ .
If instead one adopts the law $`v_\gamma =[dE/dp]_{m=0}`$, still assuming (19), one obtains
$$v_\gamma (\nu )=\frac{dE}{dp}=\frac{\mathrm{cosh}^{\frac{3}{2}}\frac{\sqrt{2}E}{E_p}}{\mathrm{cosh}\frac{E}{\sqrt{2}E_p}}$$
(71)
and then the flux density takes the form
$$I_\nu =16\pi AE_p^2\nu \frac{\mathrm{sinh}^2\frac{\nu }{\sqrt{2}E_p}}{\left(e^{\frac{\nu }{T}}1\right)\mathrm{cosh}\frac{\sqrt{2}\nu }{E_p}}.$$
(72)
From this one easily verifies that the effects induced by the Planck-scale deformation in the case $`v_\gamma =[dE/dp]_{m=0}`$ are essentially of the same type encountered in the cases $`v_\gamma =1`$ and $`v_\gamma =p/E`$, but the quantitative differences between the case $`v_\gamma =[dE/dp]_{m=0}`$ and the other two cases are more significant then the ones between the cases $`v_\gamma =1`$ and $`v_\gamma =p/E`$. As mentioned, in absence of the Planck-scale effects the typical frequency of the photons contributing to the energy flux grows linearly with the temperature, while in the cases in which the Planck-scale effects of (19) are introduced with $`v_\gamma =1`$ or $`v_\gamma =p/E`$ the typical frequency saturates at a Planckian value. If the same Planck-scale effects are introduced with $`v_\gamma =[dE/dp]_{m=0}`$, as implicitly codified in (72), one finds that the growth of the typical frequency with temperature also slows down significantly at high temperatures but it does not completely saturate: at high temperatures the typical frequency grows logarithmically with the temperature.
In summary the choice of the speed law does not appear to affect the core features of the analysis, but it appears that it could in some cases introduce some significant quantitative differences.
## 9 Comparison with previous analyses
To our knowledge, the one we reported here, in spite of its preliminary nature, is at this point the most composite effort of exploration of the implications of a MDR and/or a GUP in black-hole thermodynamics. But parts of the overall picture we attempted to provide had been investigated previously, and it seems appropriate to comment briefly on this previous related studies.
Closest in spirit to our perspective are the studies of the implications of the GUP for black-hole thermodynamics reported in Refs. and . Whereas for us (2) is to be handled prudently, as it could possibly be only an approximate form of a more complicated all-order-in-$`L_p`$ formula, in Refs. the formula (2) is taken as the exact form of the GUP, thereby leading to a corresponding form of the entropy-area relation. Perhaps more importantly Refs. assume that the GUP would not affect the black-body spectrum and in particular a standard expression for Stefanโs law is used even in Planckian regimes. There was no investigation of MDRs in Refs. .
An attempt to describe Hawking radiation in presence of a MDR was reported in Ref. . There the problem is approached from the field-theoretic perspective, considering possible modification of the field equations coming from the MDR. No explicit formula for the corrections to the Hawking spectrum and to the entropy-area relation was obtained in Ref. .
Ref. investigates how a general form of the GUP could modify the volume element of phase space, and therefore the black-body-radiation formula, using the Hamiltonian formulation in terms of Poisson brackets.
In Ref. an analysis of black-body radiation is carried out in presence of a MDR of the type emerging from a proposed โsemiclassical limitโ of Loop Quantum Gravity, which is analogous to the โleading orderโ MDR (1) we studied in some parts of this paper. The results reported there are consistent with the power-series formulas for Stefanโs and Wienโs law which we derived. The features we exposed in considering some illustrative examples of all-order MDRs, were not discussed in Ref. . Also the entropy-area relation and the aspects of black-hole evaporation which we considered here were not part of the analysis reported in Ref. , and Ref. did not consider the possibility of a GUP.
Ref. is closest in spirit to the part of our analysis where we focused on the black-body radiation spectrum, as affected by a MDR. Although the formal setup differs in several points, the results are roughly consistent with ours, including the possibility of โsaturationโ of the characteristic frequency at $`TE_p`$. There was however no investigation of the entropy-area relation and the Generalized Second law in Ref. , and Ref. also did not consider the possibility of a GUP.
## 10 Outlook
The technical difficulties that are encountered in most approaches to the quantum-gravity problem usually only allow one to grasp a few disconnected aspects of the physical picture that the theories could provide. And in some approaches even the few โphysicalโ results that are obtained, are only derived within approximation schemes whose reliability is not fully established. We have argued that in this situation it might be particularly valuable to establish a few logical links connecting some apparently unrelated aspects of the physical picture. And we showed that such a link can be found between some aspects of quantum-gravity research which have attracted strong interest in recent times, a link providing a connection between results on modified energy-momentum dispersion relations and/or modified position-momentum uncertainty principles and results on the thermodynamics of black holes. We have provided a description of log corrections to the entropy-area law for black holes that is based on the availability of a MDR and/or a GUP.
In exploring other aspects of black-hole thermodynamics as affected by MDRs and GUPs we stumbled upon a few noteworthy points. We found that the Generalized Second Law of thermodynamics might be robust enough to survive the introduction of these Planck-scale effects. We found that a MDR introducing a minimum value for wavelengths (even when no maximum value for frequencies is introduced) could lead to a description of black-body radiation in which the characteristic frequency of the radiation never exceeds a finite Planckian value (described in Eq. (50)). This in turn also affects black-hole evaporation in such a way that the temperature diverges already when the mass of the black hole decreases to a Planck-scale value (instead of diverging only in the zero-mass limit as usually assumed).
A key test for our line of analysis will come from future improved analyses within the loop-quantum-gravity approach. According to the perspective we adopted some preliminary results on the emergence of modifications of the dispersion relation that depend linearly on the Planck length (at low energies) would be incompatible with the loop-quantum-gravity results on log corrections to the entropy-area relation for black holes. We predict that improved analyses of the loop-quantum-gravity approach should lead to the emergence of a picture that is instead compatible with the conceptual link we are proposing.
As stressed in Section 8 one aspect of our analysis in which we took a rather conservative attitude (in comparison with the possibilities considered in the literature) is the one concerning the description of the speed of photons, which we assumed to be still frequency independent. We do not expect major obstacles for a generalization of our analysis with the inclusion of the possibility of a frequency-dependent speed of photons, and the preliminary investigation reported in Section 8 suggests that some of the core features that emerged from our analysis are only moderately affected by the choice of law for the speed of photons.
## Acknowledgments
The work of M. A. was supported by a Fellowship from The Graduate School of The University of North Carolina. M. A. also thanks the Department of Physics of the University of Rome for hospitality. Y. L. is partly supported by NSFC (No.10205002,10405027) and SRF for ROCS,SEM.
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# Higher genus Icosahedral Painlevรฉ curves
## 1. Introduction
A Painlevรฉ curve $`\mathrm{\Pi }`$ is an algebraic curve supporting a solution to Painlevรฉโs sixth equation (henceforth $`\text{P}_{\text{VI}}`$). That is, there should be rational functions $`y,t`$ on $`\mathrm{\Pi }`$ such that
(1)
$$t:\mathrm{\Pi }^1$$
is a Belyi map (so expresses $`\mathrm{\Pi }`$ as a branched cover, ramified only over $`0,1,\mathrm{}`$) and $`y`$ (viewed as a function of $`t`$) solves a $`\text{P}_{\text{VI}}`$ equation.
This notion was introduced by Hitchin who found an infinite family of examples related to the Poncelet problem. In essence he showed that all the modular curves $`X_1(n)`$ are Painlevรฉ curves, at least for $`n`$ prime. More precisely one should first pull back along the standard map $`X(2)X(1)`$ (with Galois group $`\text{Sym}_3=\text{PSL}_2(2)`$), so there is a diagram:
$$\begin{array}{cccc}& \mathrm{\Pi }& & X_1(n)\\ & t& & \\ ^1& X(2)& & X(1).\end{array}$$
In particular, for $`n=5`$, Hitchin wrote down the first explicit genus one Painlevรฉ curve.
The aim of this article is to write down some other explicit Painlevรฉ curves not in the above family of examples.
The (nonlinear) $`\text{P}_{\text{VI}}`$ equation controls the โisomonodromicโ (or monodromy preserving) deformations of (linear) rank two Fuchsian systems on $`^1`$ with four singularities, at $`0,t,1,\mathrm{}`$. The monodromy of such a system is a representation
$$\rho :_3=\pi _1(^1\{0,t,1,\mathrm{}\})\text{SL}_2()$$
and one of the main properties of these Painlevรฉ curves is that the monodromy of the cover (1), i.e. its permutation representation $`_2\{1,2,\mathrm{},\mathrm{deg}(t)\}`$, coincides with the standard action of the pure mapping class group of the four-punctured sphere ($`_2`$) on the orbit it generates through the conjugacy class of the representation $`\rho `$.
Hitchinโs examples arose by seeking such isomonodromic deformations when the image of $`\rho `$ was equal to a binary dihedral group, and in a previous article the author studied the case when the monodromy group is equal to the binary icosahedral group. All such solutions were classified and explicit formulae were written down for all but $`8`$ of the $`52`$ cases, including all those of genus zero and most of the genus one cases. (Five interesting cases had previously appeared in .)
Unfortunately the icosahedral Painlevรฉ curves of genus $`2`$ were not amenable to the method of construction used in , essentially due to the large degrees of the Belyi maps $`t`$. (The method used was to first obtain, from the icosahedral linear monodromy, the precise asymptotics of the $`\text{P}_{\text{VI}}`$ solution, using (the authorโs correction of) Jimboโs asymptotic formula; this determined the Puiseux expansions to arbitrary order which in turn enabled the curve to be obtained algebraically.)
However it turns out that there is a trick to convert earlier icosahedral Painlevรฉ curves (that were found in , or were previously known) into those of higher genus. Namely one may use the so-called โquadratic transformationsโ introduced by Kitaev in 1991 and written in simpler form by Ramani et al. (we learnt of them from the recent article ). Somewhat miraculously the solutions that can be obtained in this way are almost exactly the complement of those we were able to obtain by the previous method (there is a small overlap though).
Thus our aim is to explain how the quadratic transformations may be applied in this way and write down the resulting curves. (This is not entirely trivial since, if applied blindly, the quadratic transformations lead to badly parameterised solutions, for example with the wrong genus.) We also make some effort to obtain nice models (over $``$) of the resulting Painlevรฉ curves.
For example the following result will be established:
###### Theorem.
There are precisely two non-hyperelliptic icosahedral Painlevรฉ curves. The first supports two inequivalent Painlevรฉ solutions and is of genus three and isomorphic to the smooth plane quartic with affine equation
$$5(p^4+q^4)+6(p^2q^2+p^2+q^2)+1=0.$$
The second is of genus seven and is birationally isomorphic over $``$ to the affine curve cut out by the octic
$$9(p^6q^2+p^2q^6)+18p^4q^4+4(p^6+q^6)+26(p^4q^2+p^2q^4)+8(p^4+q^4)+57p^2q^2+20(p^2+q^2)+16$$
whose closure in $`^2`$ only has double point singularities. Moreover the obvious symmetries of these curves (negating and exchanging $`p`$ and $`q`$, generating a dihedral group of order $`8`$) correspond to the Okamoto symmetries of the Painlevรฉ solutions.
## 2. Background
We will constrain ourselves to giving the notation and terminology that we will use, referring the reader to or the review article and references therein for more details and geometrical background.
The sixth Painlevรฉ equation ($`\text{P}_{\text{VI}}`$) is:
$`{\displaystyle \frac{d^2y}{dt^2}}=`$ $`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{1}{y}}+{\displaystyle \frac{1}{y1}}+{\displaystyle \frac{1}{yt}}\right)\left({\displaystyle \frac{dy}{dt}}\right)^2\left({\displaystyle \frac{1}{t}}+{\displaystyle \frac{1}{t1}}+{\displaystyle \frac{1}{yt}}\right){\displaystyle \frac{dy}{dt}}`$
$`+{\displaystyle \frac{y(y1)(yt)}{2t^2(t1)^2}}\left((\theta _41)^2{\displaystyle \frac{\theta _1^2t}{y^2}}+{\displaystyle \frac{\theta _3^2(t1)}{(y1)^2}}+{\displaystyle \frac{(1\theta _2^2)t(t1)}{(yt)^2}}\right)`$
where $`\theta =(\theta _1,\theta _2,\theta _3,\theta _4)`$ are (complex) constants. This arises naturally when one tries to isomonodromically deform Fuchsian systems of the form
(2)
$$\frac{d}{dz}\left(\frac{A_1}{z}+\frac{A_2}{zt}+\frac{A_3}{z1}\right),A_i๐ค:=๐ฐ๐ฉ_2()$$
as the second pole position $`t`$ varies in $`^1\{0,1,\mathrm{}\}`$. (The parameters $`\theta `$ specify the eigenvalues of the residues: namely $`A_i`$ has eigenvalues $`\pm \theta _i/2`$ for $`i=1,2,3,4`$, where $`A_4=_1^3A_i`$.) Geometrically $`\text{P}_{\text{VI}}`$ can (thus) be thought of as the explicit form of the simplest nonabelian GaussโManin connection.
###### Definition 1.
An algebraic solution of $`\text{P}_{\text{VI}}`$ consists of a triple $`(\mathrm{\Pi },y,t)`$ where $`\mathrm{\Pi }`$ is a compact (possibly singular) algebraic curve and $`y,t`$ are rational functions on $`\mathrm{\Pi }`$ such that:
$``$ $`t:\mathrm{\Pi }^1`$ is a Belyi map (i.e. $`t`$ expresses $`\mathrm{\Pi }`$ as a branched cover of $`^1`$ which only ramifies over $`0,1,\mathrm{}`$), and
$``$ Using $`t`$ as a local coordinate on $`\mathrm{\Pi }`$ away from ramification points, $`y(t)`$ should solve $`\text{P}_{\text{VI}}`$, for some value of the parameters $`\theta `$.
Indeed given an algebraic solution in the form of a polynomial relation $`F(y,t)=0`$ one may take $`\mathrm{\Pi }`$ to be the closure in $`^2`$ of the affine plane curve defined by $`F`$. That $`t`$ is a Belyi map on $`\mathrm{\Pi }`$ follows from the Painlevรฉ property of $`\text{P}_{\text{VI}}`$: solutions will only branch at $`t=0,1,\mathrm{}`$ and all other singularities are just poles. The reason we prefer this reformulation is that often the polynomial $`F`$ is quite complicated and usually there are much simpler models of the plane curve defined by $`F`$. (The polynomial $`F`$ can of course be recovered as the minimal polynomial of $`y`$ over $`(t)`$.)
We will say a Painlevรฉ curve $`\mathrm{\Pi }`$ is โminimalโ or an โefficient parameterisationโ if $`y`$ generates the field of rational functions on $`\mathrm{\Pi }`$, over $`(t)`$, so that $`y`$ and $`t`$ are not pulled back from another curve covered by $`\mathrm{\Pi }`$ (i.e. that $`\mathrm{\Pi }`$ is birational to the curve defined by $`F`$).
The main invariants of an algebraic solution are the genus of a (minimal) Painlevรฉ curve $`\mathrm{\Pi }`$ and the degree of the corresponding Belyi map $`t`$ (the number of branches the solution has over the $`t`$-line).
We will say that two solutions of $`\text{P}_{\text{VI}}`$ are equivalent if they are related by Okamotoโs affine $`F_4`$ Weyl group symmetries of $`\text{P}_{\text{VI}}`$ (which act on the set of parameters $`\{\theta \}^4`$ in the standard way). (See e.g. for formulae for this action.) For an algebraic solution, this acts within the set of rational functions on the curve $`\mathrm{\Pi }`$, and preserves the degree and genus of the solution (at least if the linear monodromy representation is irreducible and not rigid).
We are interested here in the case where the monodromy group of the linear system (2) is equal to the binary icosahedral group<sup>1</sup><sup>1</sup>1more precisely we are interested in the solutions equivalent to such; one should bear in mind that the Okamoto transformations can change the monodromy group, and it will in fact be simpler to work at different equivalent values of the parameters $`\theta `$. cf. Remark 4 $`\mathrm{\Gamma }\text{SL}_2()`$. To understand the different cases that may occur essentially amounts to studying the different conjugacy classes of the local projective monodromies. Recall that the icosahedral rotation group $`\mathrm{\Gamma }/\pm A_5\text{SO}_3()`$ has four non-trivial conjugacy classes, which we will label $`a,b,c,d`$ corresponding to rotations by $`\frac{1}{2},\frac{1}{3},\frac{1}{5},\frac{2}{5}`$-of a turn, respectively. Thus we define, as in , the $`A_5`$-type of a representation
$$\rho :\pi _1(^1\{0,t,1,\mathrm{}\})\mathrm{\Gamma }$$
to be the corresponding unordered set of four conjugacy classes of projective local monodromies (i.e. take the conjugacy classes of the images in $`A_5`$ of the elements $`\rho (\gamma _i)`$ for simple loops $`\gamma _i`$ encircling one of $`0,t,1`$ or $`\mathrm{}`$ once). The different cases that occur are tabulated in .
Two inequivalent icosahedral solutions will be said to be siblings if their monodromy representations $`\rho `$ are related by the nontrivial outer automorphism of $`A_5`$ (swapping the conjugacy classes $`c,d`$). They will have the same Belyi map $`t`$, just a different solution function $`y`$. (In general it is useful to generalize this notion by considering Galois conjugate representations, e.g. for representations into the $`237`$ triangle group there are sometimes three siblings, cf. .)
## 3. Quadratic transformations
The basic idea behind the quadratic transformations is as follows. Given an icosahedral Fuchsian system $`A`$ with $`A_5`$ type $`a^2\xi \eta `$ for some $`\xi ,\eta \{a,b,c,d\}`$ (i.e. with two local monodromies, say at $`0`$ and $`\mathrm{}`$, of order two in $`\text{PSL}_2()`$) we can pull back along the map $`wz=w^2`$ to get a Fuchsian system with two apparent singularities at $`0`$, $`\mathrm{}`$ and four non-apparent singularities at $`\pm 1,\pm \sqrt{t}`$. Removing the apparent singularities (using Schlesinger transformations) yields a system $`B`$ with $`A_5`$ type $`\xi ^2\eta ^2`$, which may be put in the form (2) by a coordinate transformation. Isomonodromic deformations of $`A`$ correspond to isomonodromic deformations of $`B`$, and one can obtain formulae relating the corresponding $`\text{P}_{\text{VI}}`$ solutions. In practice the formulae are much simpler at different (Okamoto equivalent) values of the parameters (see Ramani et al. (2.7)). We should emphasise that these transformations are not really symmetries of the family of Painlevรฉ VI equations since the conditions on the parameters restrict us to a co-dimension two subset of the four-dimensional parameter space. Nonetheless they are precisely what is needed to obtain the eight outstanding icosahedral solutions, since they all have the desired factor of $`a^2`$ in their $`A_5`$ types. Indeed for these cases, this procedure gives an algebraic relation with a solution having half the number of branches; Examining table 1 of we see solution 31 $``$ solution 44 and in turn solution 44 $``$ solution 50. Similarly
$$324551,3947,4048,414952.$$
The formula of Ramani et al. that we will use to construct these outstanding solutions from known solutions is as follows. (In fact this is the inverse of the formula (2.7), having converted their parameters to our conventions.)
###### Proposition 2 ().
Given a solution $`(y_0,t_0)`$ of $`\text{P}_{\text{VI}}`$ with parameters of the form $`\theta =(0,\theta _2,\theta _3,1)`$ then, by taking two square roots, one obtains a new solution $`(y,t)`$ with parameters $`\theta =(\theta _3,\theta _2,\theta _2,2\theta _3)/2`$ where
$$y=\frac{(\tau 1)(\eta +1)}{(\tau +1)(\eta 1)},t=\left(\frac{\tau 1}{\tau +1}\right)^2$$
with $`\eta ^2=y_0,\tau ^2=t_0`$.
Note that negating $`\tau `$ corresponds to the Okamoto symmetry $`(y,t)(y/t,1/t)`$ and negating both $`\eta `$ and $`\tau `$ corresponds to $`(y,t)(1/y,1/t).`$
In practice this will usually lead to an inefficiently parameterised Painlevรฉ curve. In the cases at hand this may be remedied as follows. (In the process we will convert the formula to that most directly useful to us.) The relation between the Painlevรฉ curve $`\mathrm{\Pi }^{}`$ we end up with and the original curve $`\mathrm{\Pi }`$ may be summarised by the diagram:
where the numbers indicate the degrees of the maps, and $`\stackrel{~}{\mathrm{\Pi }}`$ is the intermediate curve obtained by adjoining the two square roots to the function field of $`\mathrm{\Pi }`$.
Suppose our initial solution is a pair of functions of the form
(3)
$$Y=\frac{1}{2}+a_Y(s)u,T=\frac{1}{2}+a_T(s)u$$
for parameters of the form $`\theta =(0,\theta _2,0,\theta _4)`$ on a curve of the form
$$\mathrm{\Pi }:=\{u^2=u_2(s)\}$$
where $`u_2`$ is a polynomial, and $`a_Y,a_T`$ are rational functions of $`s`$. In other words $`\mathrm{\Pi }`$ is a double cover of the $`s`$-line $`_s^1`$, and the symmetry of $`\mathrm{\Pi }`$ (negating $`u`$) corresponds to the symmetry $`(y,t)(1y,1t)`$. Our basic observation is that the parameter $`u`$ will drop out in the solution obtained, as follows.
Applying the Okamoto transformation $`(Y,T)(Y/(Y1),T/(T1))`$ yields a solution to which we may apply Proposition 2. Thus we need to take square roots of $`Y/(Y1),T/(T1)`$, i.e. of expressions of the form $`(A+u)/(Au)`$ where $`A=2au_2`$ is still a rational function of $`s`$. A useful trick is to look for square roots of similar form: i.e. to find $`B`$ such that
$$\left(\frac{B+u}{Bu}\right)^2=\frac{A+u}{Au}.$$
Taking the square root of both sides and solving for $`B`$ we find
$$B=A\pm \sqrt{A^2u_2}$$
which does not involve $`u`$. Carrying this out for both $`Y`$ and $`T`$ we obtain
$$\eta =\frac{B_Y+u}{B_Yu},\tau =\frac{B_T+u}{B_Tu},$$
where $`B_i=A_i\pm \sqrt{A_i^2u_2}`$ for $`i=Y,T`$. Then the formulae of Proposition 2 yield
$$y=\frac{B_Y}{B_T},t=\frac{u_2}{B_T^2}$$
neither of which involves $`u`$. Thus $`\mathrm{\Pi }^{}`$ can be viewed as either the quotient of $`\stackrel{~}{\mathrm{\Pi }}`$ by the involution negating $`u`$ or as the four-fold cover of the $`s`$-line obtained by adjoining functions $`v,w`$ with
$$v^2=A_Y^2u_2,w^2=A_T^2u_2$$
where $`A_i=2u_2a_i`$ for $`i=Y,T`$. (The reader may verify that the involution of $`\mathrm{\Pi }^{}`$ negating both $`v`$ and $`w`$ together yields the transformation $`(y,t)(1/y,1/t)`$.)
In turn if we apply the transformation $`(y,t)(y/(y1),t/(t1))`$ we will obtain a solution of form similar to (3). In summary (after some relabelling) the version of the quadratic transformations we will actually use is as follows:
###### Corollary 3.
If the functions $`y_0,t_0`$ of the form
$$y_0=\frac{1}{2}+a_y(s)u,t_0=\frac{1}{2}+a_t(s)u$$
are a $`\text{P}_{\text{VI}}`$ solution with parameters $`\theta =(0,\theta _2,0,\theta _4)`$ on a Painlevรฉ curve of the form
$$\mathrm{\Pi }:=\{u^2=u_2(s)\}$$
for a polynomial $`u_2(s)`$, then the functions
$$y=\frac{1}{2}+\frac{w+v}{2(A_yA_t)},t=\frac{1}{2}\frac{A_t}{2w}$$
are a $`\text{P}_{\text{VI}}`$ solution for parameters $`\theta =(1\theta _4,\theta _2,1\theta _4,2\theta _2)/2`$ on the curve obtained by adjoining to $`(s)`$ the functions $`v,w`$ where
$$v^2=A_y^2u_2,w^2=A_t^2u_2$$
and $`A_i=2a_iu_2`$ for $`i=y,t`$.
Of course, a similar result is true upon replacing $`_s^1`$ by an arbitrary genus curve, but this will be sufficient for us here. Note that negating both $`v`$ and $`w`$ now corresponds to the Okamoto transformation $`(y,t)(1y,1t)`$.
## 4. Solutions
We will now carry out the following steps to find the formulae for the outstanding icosahedral solutions:
1) Choose an icosahedral solution from the table in and if possible convert it, via Okamoto transformations, into a solution with parameters of the form $`(0,\theta _2,0,\theta _4)`$,
2) Apply Corollary 3 to obtain a new solution, which will (in the examples here) have twice the number of branches (and larger genus) than the original solution,
3) Look for a simple model of the resulting Painlevรฉ curve (either as a double cover of some $`^1`$, if it is hyperelliptic, or as a low degree plane curve otherwise).
A priori suitable solutions for step 1) are easily detected by looking for two zero coordinates in the solutionโs alcove point listed in table 1 of .
###### Remark 4.
To aid the interested reader, and avoid typos, a Maple text file of the solutions of this article has been included with the source file (obtained by clicking on โOther formatsโ) for the preprint version on the math arxiv. This file also contains solutions equivalent to those written here for which the corresponding isomonodromic family of Fuchsian systems has finite (icosahedral) monodromy group.
10 branch genus zero $``$ 20 branch genus one.
Applying some Okamoto transformations to the $`H_3`$ solution from E.33, , which is equivalent to icosahedral solution $`32`$, one obtains the solution
$$y_0=\frac{1}{2}\frac{\left(3s^2+6s1\right)u}{16s^2},t_0=\frac{1}{2}+\frac{uP}{256\left(5s1\right)s^3}$$
for parameters $`\theta =(0,1/5,0,1)`$ where $`u^2=s`$ and $`P=27s^5315s^4370s^3+170s^225s+1.`$ Applying Corollary 3 to this (and adjusting $`v,w`$ slightly to remove square factors) yields the solution
$$y=\frac{1}{2}\frac{16s\left(5s1\right)+vw}{2\left(s1\right)\left(3s+1\right)v},t=\frac{1}{2}\frac{P}{2\left(s1\right)v^2w}$$
for parameters $`\theta =(0,1,0,9)/10`$ with $`P`$ as above and where $`w=vw_1`$ and
(4)
$$v^2=(9s1)(s1),w_1^2=s^218s+1.$$
One may check directly that this is a genus one solution with twenty branches, and is equivalent to icosahedral solution $`45`$. (It is reassuring to compute the monodromy of the cover $`t:\mathrm{\Pi }^{}^1`$ and find it has the properties listed in table 1 of .) Our next aim is to find a good model of the elliptic curve defined by (4), preserving the symmetry negating $`v`$. We will do this by parameterising the conic $`w_1^2=s^218s+1`$ as follows:
$$s=\frac{j^21}{2j18},w_1=\frac{j^218+1}{2j18}.$$
Then if we define $`v=\frac{z}{2j18}`$ the condition that $`v^2=(9s1)(s1)`$ says that $`(z,j)`$ is a point of the elliptic curve
(5)
$$z^2=(9j^22j+9)(j^22j+17),$$
and the above formulae give $`y,t`$ explicitly as functions on this curve. (One may show, using Magma for example, that this elliptic curve corresponds to entry 200B1 of Cremonaโs tables of elliptic curves and for example is isomorphic over $``$ to the plane cubic $`u^2=s(s^25s+5)`$, but this model hides the symmetry of the Painlevรฉ solution.)
Similarly we can proceed with the sibling solution to that above, to obtain the solution:
$$y=\frac{1}{2}\frac{64\left(5s1\right)s^2+\left(s1\right)vw}{2\left(3s^3+75s^215s+1\right)v}$$
with $`t,s,v,w`$ as above but $`\theta =(0,3,0,7)/10`$. This is equivalent to icosahedral solution $`44`$.
20 branch genus one $``$ 40 branch genus three.
We can apply Corollary 3 again to the resulting solutions above, since their parameters are again of the desired form. Solution $`45`$ then yields the solution
$$y=\frac{1}{2}+\frac{\left(j^218j+1\right)z^2+16\left(j+3\right)\left(j+1\right)vw}{8\left(3j7\right)\left(j9\right)\left(j1\right)^2v},t=\frac{1}{2}+\frac{uP}{256\left(5s1\right)s^3}$$
with $`\theta =(1,1,1,19)/20`$, where $`P(s),z^2`$ are as in the previous subsection,
$$s=\frac{j^21}{2j18},u=\frac{w}{2j18}$$
and now
(6)
$$v^2=(j1)(j9)(5j^22j+13),w^2=2(j9)(j^21).$$
One may check directly that this is a genus 3 solution with forty branches and is equivalent to icosahedral solution $`51`$. (Note that $`t`$ is simply the pullback of the original degree $`10`$ function $`t_0`$.) The curve defined by (6) is not hyperelliptic, so we can find a plane model by taking the canonical embedding. (Eliminating $`s`$ from the equations (6) yields a singular plane sextic, and we compute three independent differentials directly on this.) This gives the following model of the Painlevรฉ curve as a smooth plane quartic, with affine equation:
(7)
$$5(p^4+q^4)+6(p^2q^2+p^2+q^2)+1=0.$$
The solution functions $`(y,t)`$ become functions on this quartic by setting
$$v=\frac{200p(6p^2+5q^2+1)}{84p^2q^255q^4166q^2156p^231},w=qv,s=\frac{28v^24w^2+800}{3v^2+15w^2800}.$$
Notice that this curve has three involutions (generating a group isomorphic to the dihedral group of order eight). These correspond to the Okamoto symmetries coming from the three hyperplanes on which the solutionโs parameters lie (as listed in table 1 of ). In more detail the symmetries mapping $`(p,q)`$ to $`(p,q),(p,q),(q,p)`$ correspond to the Okamoto symmetries mapping $`(y,t)`$ to
$$(1y,1t),(\frac{y(t1)}{ty},1t),(\frac{yt}{y1},t)$$
respectively.
Similarly, from the sibling solution 44 one obtains the following, which is equivalent to icosahedral solution $`50`$:
$$y=\frac{1}{2}+\frac{\left(j^218j+1\right)\left(j^22j+17\right)z^2+8\left(j1\right)\left(j^3+57j^269j+75\right)vw}{8\left(3j^321j^215j31\right)w^2v}$$
with $`t,v,w,z^2`$ as above and $`\theta =(3,3,3,17)/20`$.
15 branch genus one $``$ 30 branch genus two.
If we apply some Okamoto transformations to icosahedral solution 39 (from ) then we can obtain the solution
$$y_0=\frac{1}{2}\frac{u\left(2s^2+3s3\right)}{6\left(s+1\right)\left(4s^2+15s+15\right)},t_0=\frac{1}{2}\frac{uP}{18\left(4s^2+15s+15\right)^2\left(s^25\right)}$$
with $`\theta =(0,7/15,0,13/15)`$, where $`u^2=3(s+5)(4s^2+15s+15)`$ and
$$P=2s^7+10s^690s^4135s^3+297s^2+945s+675.$$
Applying Corollary 3 to this, and again adjusting $`v,w`$ to remove square factors, yields the solution:
$$y=\frac{1}{2}+\frac{\left(s^25\right)u^2v+s\left(s3\right)\left(s+1\right)w^3}{2\left(s3\right)\left(s+5\right)\left(s^3+s^29s15\right)w^2},t=\frac{1}{2}+\frac{\left(s+5\right)^2P}{4s\left(s^29\right)w^3}$$
with $`\theta =(2,7,2,23)/30`$ where $`P`$ and $`u^2`$ are as above and
(8)
$$v^2=s\left(s+5\right)\left(s+2\right)\left(s3\right),w^2=s\left(s+5\right)\left(s+2\right)\left(s+3\right).$$
This has thirty branches, genus two and is equivalent to icosahedral solution $`47`$. Being of genus two, the curve (8) is hyperelliptic. We will express it as a double cover of a $`^1`$ branched at six points. Indeed by choosing a parameter $`j`$ on the conic $`x^2=s^29`$, the Painlevรฉ curve (8) becomes isomorphic to the hyperelliptic curve
$$z^2=(j^2+9)(j+9)(j+1)(j^2+4j+9)$$
via the map
$$v=\frac{j3}{4j^2}z,w=\frac{j+3}{4j^2}z,s=\frac{j^2+9}{2j}.$$
Similarly we can repeat starting with solution 40 (the sibling of 39) and obtain solution 48 (the sibling of solution 47). The result is
$$y=\frac{1}{2}+\frac{\left(s^25\right)u^2+\left(s^26s15\right)vw}{2s\left(s+5\right)\left(s+3\right)^2v}$$
with $`\theta =(4,1,4,29)/30`$ and with $`t,v,w,u^2,s`$ as for solution 47 above.
18 branch genus one $``$ 36 branch genus three.
Next we will start with icosahedral solution 41 (from ; the 10 page implicit form of this solution in the preprint version of is not useful here). Applying some Okamoto transformations yields the solution:
$$y_0=\frac{1}{2}\frac{8s^312s^2+3s4}{6u},t_0=\frac{1}{2}+\frac{P}{54s\left(s1\right)u^3}.$$
for $`\theta =(0,1/3,0,1)`$ where
(9)
$$u^2=s(8s^211s+8),$$
and
$$P=\left(s+1\right)\left(32(s^8+1)320(s^7+s)+1112(s^6+s^2)2420(s^5+s^3)+3167s^4\right).$$
Applying Corollary 3 to this yields the solution:
$$y=\frac{1}{2}\frac{9s\left(s1\right)u^2+\left(s2\right)wv}{2\left(s^3+12s^212s+4\right)\left(2s1\right)v},t=\frac{1}{2}\frac{P}{4\left(2s1\right)\left(s2\right)v^2w}$$
with $`P,u^2`$ as above, $`\theta =(0,1,0,5)/6`$ and $`w=vw_1`$ where
(10)
$$v^2=(s2)(2s1)(2s^2+s+2),w_1^2=s^27s+1.$$
One may check this is a $`36`$ branch genus three solution and is equivalent to icosahedral solution $`49`$. However in this case the curve defined by (10) is hyperelliptic. Indeed let $`j`$ be a parameter on the conic $`w_1^2=s^27s+1`$, so for example
(11)
$$w_1=\frac{j^27j+1}{2j7},s=\frac{j^21}{2j7}.$$
Then the Painlevรฉ curve (10) becomes isomorphic to
$$z^2=\left(j^24j+13\right)\left(2j^22j+5\right)\left(2j^4+2j^33j^258j+107\right)$$
via (11) and the assignment $`v=z/(2j7)^2`$.
36 branch genus three $``$ 72 branch genus seven.
Finally we can apply Corollary 3 again to the solution above (since the parameters are of the desired form) to obtain the largest icosahedral solution. The solution is given by
$$y=\frac{1}{2}+\frac{9\left(j1\right)\left(j^3+27j^257j+79\right)wv+2\left(2j^22j+5\right)\left(j^27j+1\right)\left(2j^4+2j^33j^258j+107\right)\left(j^24j+13\right)^2}{6\left(j^21\right)\left(2j^2+j+17\right)\left(j^33j^2+3j11\right)\left(2j7\right)^2v},$$
$$t=\frac{1}{2}+\frac{P}{54s\left(s1\right)u^3},$$
where $`P(s)`$ is the polynomial in the previous subsection,
$$s=\frac{j^21}{2j7},u=\frac{w}{(2j7)^2},$$
and
(12)
$$v^2=\left(j+1\right)\left(6+j^22j\right)\left(4j^213j+19\right),$$
(13)
$$w^2=\left(j1\right)\left(2j7\right)\left(j+1\right)\left(2j^2+j+17\right)\left(4j^213j+19\right).$$
Note that equation (13) is equivalent to equation (9) so $`t`$ is the pullback of the original degree $`18`$ function $`t_0`$.
One may check directly that this does indeed define a genus seven, $`72`$ branch Painlevรฉ solution and is equivalent to icosahedral solution $`52`$. Of course being of genus $`7`$ the degreeโgenus formula implies we cannot hope to find a non-singular plane model of the Painlevรฉ curve. Instead we will look for a low degree plane model with mild singularities. (The curve obtained upon eliminating $`j`$ from (12), (13) is a highly singular degree $`14`$ plane curve, with large coefficients.) We do this by selecting a three-dimensional subspace of the space of holomorphic one-forms on the curve, and taking the corresponding plane curve. After some trial and error choosing a good subspace we found the following plane octic with only double points (ten nodes and two tacnodes):
$$9(p^6q^2+p^2q^6)+18p^4q^4+4(p^6+q^6)+26(p^4q^2+p^2q^4)+8(p^4+q^4)+57p^2q^2+20(p^2+q^2)+16.$$
The map between the curves is given by
$$p=\frac{w}{3\left(j1\right)v},q=\frac{v}{3(j^22j+6)}$$
and, if needed, the (rather long) inverse appears in the accompanying computer file (see Remark 4).
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# Finite-temperature properties of hard-core bosons confined on one-dimensional optical lattices
## I Introduction
The study of ultracold quantum gases loaded on optical lattices has become a very active area of experimental and theoretical reseach in recent years. Optical lattices enable enhancing interactions between atoms in weakly interacting Bose-Einstein condensates (BECโs) and reducing the effective dimensionality of the system. They allow the experimental realization of strongly correlated bosons, well described by the Bose-Hubbard Hamiltonian fisher89 ; jaksch98 , with the consequent observation of the superfluidโMott-insulator transition greiner02 ; wessel04 . In addition, optical lattices have been used to obtain one-dimensional (1D) systems moritz03 ; tolra04 and to examine the superfluidโMott-insulator transition in 1D stoferle03 ; batrouni02 .
Due to the strong effects of quantum fluctuations and the possibility of obtaining exact theoretical results, 1D systems are a very attractive laboratory for both experiments and theory. Theoretically, it was shown by Olshanii that in 1D in regimes of large scattering length, low densities, and low temperatures bosons behave as a gas of impenetrable particles known as hard-core bosons (HCBโs) olshanii98 . Such a 1D gas (also recently called a Tonks-Girardeau gas) was introduced by Girardeau, who established an exact mapping between these strongly correlated bosons and noninteracting spinless fermions girardeau60 . Since then 1D HCBโs have been extensively studied by different techniques in both homogeneous lieb63 ; lenard64 ; vaidya79 ; haldane81 ; korepin93 ; cazalilla04\_1 and harmonically trapped girardeau01 ; forrester03 ; gangardt04 systems.
The experimental realization of 1D HCBโs followed after more than 40 years of the theoretical introduction of the model paredes04 ; kinoshita04 , with paredes04 and without kinoshita04 an additional lattice along the 1D axis. The additional 1D lattice paredes04 facilitates the achievement of the HCB regime with respect to the continuum case. It allows experimentalists to change the effective mass of the particles and, consequently, the ratio between interaction and kinetic energies paredes04 . Although at very low densities (when interparticle distances are much larger than the lattice spacing) HCBโs on a lattice are equivalent to HCBโs in continuous space, this is not the case for arbitrary fillings cazalilla04 . On 1D lattices the HCB Hamiltonian can be mapped onto the 1D XY model of Lieb, Schulz, and Mattis lieb61 . For periodic systems this model has been also studied extensively in the literature mccoy68 ; vaidya78 ; mccoy83 . More recently, renewed interest has arisen on the properties of HCBโs when additional confining potentials are introduced, as the case relevant to experiments paredes04 .
Remarkably, even in trapped inhomogeneous systems power-law behavior known from the periodic case is present rigol04\_1 . The one-particle density matrix exhibits a universal power-law decay with exponent $`1/2`$ independent of the power of the confining potential rigol04\_1 . These quasi-long-range one-particle correlations generate quasicondensates with occupations scaling proportional to $`\sqrt{N_b}`$ (with $`N_b`$ the number of HCBโs in the system) rigol04\_1 . The nonequilibrium dynamics of HCBโs on 1D lattices has also been shown to display very interesting features. Quasicondensates of HCBโs emerge at finite momentum when the system starts its free evolution from a pure Mott-insulating (Fock) state rigol04\_2 . In adition, it was shown in Ref. rigol04\_3 that in 1D when there is no Mott insulator in the trap, the momentum distribution of expanding HCBโs rapidly approaches that of noninteracting fermions rigol04\_3 .
In this work we present an exact study of the finite-temperature properties of HCBโs confined on 1D optical lattices. Following the spirit of Refs. rigol04\_1 ; rigol04\_2 , we develop an exact numerical approach based on the Jordan-Wigner transformation, which maps HCBโs on a lattice onto noninteracting spinless fermions. We will focus on the effect of temperature on the off-diagonal behavior of one-particle correlations and related quantities like the momentum distribution function $`n_k`$ and the natural orbital occupations. The natural orbitals ($`\varphi ^\eta `$) are defined as the eigenfunctions of the one-particle density matrix ($`\rho _{ij}`$) penrose56 ,
$$\underset{j=1}{\overset{N}{}}\rho _{ij}\varphi _j^\eta =\lambda _\eta \varphi _i^\eta ,$$
(1)
and have occupations $`\lambda _\eta `$. (They resemble one-particle states in these strongly interacting systems.) In dilute higher-dimensional gases, when only the lowest natural orbital (the highest occupied one) scales $`N_b`$, it can be regarded as the BEC order parameter leggett01 . Here we will show that in 1D even at very low temperatures ($`T`$), when the energy of the system is almost identical to the one at $`T=0`$, the momentum distribution and the lowest natural orbital occupations can exhibit significant changes with respect to their values in the ground state.
The exposition is organized as follows. In Sec. II we describe our exact approach to study finite-temperature systems. In Sec. III we discuss the properties of HCBโs in a perfect box (an open system). HCBโs confined in harmonic traps are analyzed in Sec. IV. Since in this work we follow a grand-canonical approach to study finite-temperature properties, in Sec. V we compare exact results obtained from a grand-canonical calculation with results obtained from a canonical one for small lattice sizes, like the ones recently achieved experimentally paredes04 . Finally, the conclusions are presented in Sec. VI.
## II Exact finite-temperature approach
In this section we detail the exact approach followed to study the finite-temperature properties of HCBโs confined on 1D lattices. The HCB Hamiltonian can be written as
$$H=t\underset{i}{}\left(b_i^{}b_{i+1}+\text{H.c.}\right)+V_2\underset{i}{}x_i^2n_i,$$
(2)
with the additional on-site constraints
$$b_i^2=b_i^2=0,\{b_i,b_i^{}\}=1,$$
(3)
which avoid double or higher occupancy. The bosonic creation and annihilation operators at site $`i`$ are denoted by $`b_i^{}`$ and $`b_i`$, respectively, and the local density operator by $`n_i=b_i^{}b_i`$. The brackets in Eq. (3) apply only to on-site anticommutation relations; for $`ij`$, these operators commute as usual for bosons $`[b_i,b_j^{}]=0`$. In Eq. (2), the hopping parameter is denoted by $`t`$ and the last term represents a harmonic trap with curvature $`V_2`$.
In order to exactly calculate HCB properties, we use the Jordan-Wigner transformation jordan28
$$b_i^{}=f_i^{}\underset{\beta =1}{\overset{i1}{}}e^{i\pi f_\beta ^{}f_\beta },b_i=\underset{\beta =1}{\overset{i1}{}}e^{i\pi f_\beta ^{}f_\beta }f_i,$$
(4)
which maps the HCB Hamiltonian onto the one of noninteracting spinless fermions,
$`H_F=t{\displaystyle \underset{i}{}}\left(f_i^{}f_{i+1}+\text{H.c.}\right)+V_2{\displaystyle \underset{i}{}}x_i^2n_i^f,`$ (5)
where $`f_i^{}`$ and $`f_i`$ are the creation and annihilation operators for spinless fermions at site $`i`$ and $`n_i^f=f_i^{}f_i`$ is the local particle number operator.
The mapping as presented above is only valid for open systems, as relevant for confined bosons in experiments paredes04 ; kinoshita04 . In such cases HCBโs and fermions have exactly the same spectrum. In order to deal with 1D cyclic chains, with $`N`$ lattice sites, one needs to consider that
$$b_1^{}b_N=f_1^{}f_N\mathrm{exp}\left(i\pi \underset{\beta =1}{\overset{N}{}}n_\beta ^f\right),$$
(6)
so that when the number of particles in the system \[$`_in_i=_in_i^f=N_b`$\] is odd, the equivalent fermionic Hamiltonian satisfies periodic boundary conditions; otherwise, if $`N_b`$ is even, antiperiodic boundary conditions are required in Eq. (5).
Since for finite temperatures we will consider a grand-canonical ensembleโi.e., a system with fluctuating number of particlesโin order to avoid the dependence of the equivalent fermionic Hamiltonian on $`N_b`$ we restrict our analysis to the open case. In this case the nontrivial differences between the properties of HCBโs and fermions are only in off-diagonal correlation functions.
For finite temperatures,and within the grand-canonical formalism, the HCB one-particle density matrix can be written in terms of the equivalent fermionic system as
$`\rho _{ij}`$ $``$ $`{\displaystyle \frac{1}{Z}}\mathrm{Tr}\left\{b_i^{}b_j\mathrm{exp}\left[\left(H\mu {\displaystyle \underset{l}{}}n_l\right)/k_BT\right]\right\}`$
$`=`$ $`{\displaystyle \frac{1}{Z}}\mathrm{Tr}\{f_i^{}f_j{\displaystyle \underset{\beta =1}{\overset{j1}{}}}\mathrm{exp}(i\pi n_\beta ^f)`$
$`\times `$ $`\mathrm{exp}[(H_F\mu {\displaystyle \underset{l}{}}n_l^f)/k_BT]{\displaystyle \underset{\gamma =1}{\overset{i1}{}}}\mathrm{exp}(i\pi n_\gamma ^f)\},`$
where, in addition to Eqs. (4) and (5), we have used the cyclic property of the trace. In Eq. (II), $`\mu `$ denotes the chemical potential, $`k_B`$ the Boltzmann constant, $`T`$ the temperature of the system, and $`Z`$ the partition function
$$Z=\mathrm{Tr}\left\{\mathrm{exp}\left[\left(H_F\mu \underset{l}{}n_l^f\right)/k_BT\right]\right\}.$$
(8)
To calculate traces over the Fock space we will take advantage of the fact that in the equivalent fermionic system Fock states are Slater determinants,
$$|\mathrm{\Psi }_F=\underset{i=1}{\overset{N_f}{}}\underset{j=1}{\overset{N}{}}P_{ji}f_j^{}|0,$$
(9)
with $`N_f`$ the number of fermions and
$$๐=\left(\begin{array}{cccccc}P_{11}& P_{12}& & & & P_{1N_f}\\ P_{21}& P_{22}& & & & P_{2N_f}\\ & & & & & \\ & & & & & \\ & & & & & \\ P_{N1}& P_{N2}& & & & P_{NN_f}\end{array}\right)$$
(10)
the matrix of the components.
The action of exponentials bilinear on fermionic creation and annihilation operators, as the ones on Eqs. (II) and (8), on Slater determinants generates new Slater determinants muramatsu99 ; assaad02
$$\mathrm{exp}\left(\underset{ij}{}f_i^{}X_{ij}f_j\right)|\mathrm{\Psi }_F=\underset{i=1}{\overset{N_f}{}}\underset{j=1}{\overset{N}{}}P_{ji}^{}f_j^{}|0,$$
(11)
where
$$๐^{}=e^๐๐.$$
(12)
Using this property one can prove the following identity for the trace over the fermionic Fock space muramatsu99 ; assaad02
$`\mathrm{Tr}[\mathrm{exp}\left({\displaystyle \underset{ij}{}}f_i^{}X_{ij}f_j\right)\mathrm{exp}\left({\displaystyle \underset{kl}{}}f_k^{}Y_{kl}f_l\right)\mathrm{}`$
$`\mathrm{exp}\left({\displaystyle \underset{mn}{}}f_m^{}Z_{mn}f_n\right)]`$
$`=det\left[๐+e^๐e^๐\mathrm{}e^๐\right],`$ (13)
which immediately allows one to calculate the partition function as
$`Z`$ $`=`$ $`det\left[๐+e^{(๐_F\mu ๐)/k_BT}\right]`$ (14)
$`=`$ $`{\displaystyle \underset{i}{}}\left[1+e^{(E_{ii}\mu )/k_BT}\right],`$
where $`๐`$ is the identity matrix. The last equality was obtained after diagonalizing Hamiltonian (5), $`๐_F๐=\mathrm{๐๐}`$, $`๐`$ is the orthogonal matrix of eigenvectors, and $`๐`$ is the diagonal matrix of eigenvalues.
The trace in Eq. (II) is calculated along the same line. For $`ij`$, we notice that
$$f_i^{}f_j=\mathrm{exp}\left(\underset{mn}{}f_m^{}A_{mn}f_n\right)1,$$
(15)
where the only nonzero element of $`๐`$ is $`A_{ij}=1`$. Then, for $`ij`$, $`\rho _{ij}`$ can be obtained as
$`\rho _{ij}={\displaystyle \frac{1}{Z}}`$ $`\{det[๐+(๐+๐)๐_1๐e^{(๐\mu ๐)/k_BT}๐^{}๐_2]`$ (16)
$`det[๐+๐_1๐e^{(๐\mu ๐)/k_BT}๐^{}๐_2]\}.`$
$`๐_1`$ ($`๐_2`$) is diagonal with the first $`j1`$ ($`i1`$) elements of the diagonal equal to $`1`$ and the others equal to $`1`$.
The diagonal elements of the one-particle density matrix are the same of noninteracting fermions \[see Eq. (II) for $`i=j`$\] and can be easily calculated as muramatsu99 ; assaad02
$`\rho _{ii}`$ $`=`$ $`\left[๐+e^{(๐_F\mu ๐)/k_BT}\right]_{ii}^1`$ (17)
$`=`$ $`\left[๐\left(๐+e^{(๐\mu ๐)/k_BT}\right)^1๐^{}\right]_{ii}.`$
As usual, the chemical potential is fixed using the relation $`N_b=_i\rho _{ii}`$ to obtain the desired number of particles in the system.
## III Hard-core bosons in a box
In this section we study the finite-temperature properties of HCBโs on a perfect box. In this case the HCB Hamiltonian can be written as
$$H=t\underset{i=1}{\overset{N1}{}}\left(b_i^{}b_{i+1}+\text{H.c.}\right),$$
(18)
with the additional on-site constraints (3).
The above Hamiltonian is particle-hole symmetric, like the one of periodic systems, under the transformation $`h_i=b_i^{},h_i^{}=b_i`$ ($`h_i^{}`$ and $`h_i`$ are hole creation and annihilation operators). The particle-hole symmetry implies that the off-diagonal elements of the one-particle density matrix for $`N_b`$ HCBโs \[$`\rho _{ij}(N_b)`$\] and for $`(NN_b)`$ HCBโs \[$`\rho _{ij}(NN_b)`$\] are identical. Diagonal elements satisfy the relation $`\rho _{ii}(N_b)=1\rho _{ii}(NN_b)`$. This leads to a momentum distribution function
$$n_k=\frac{1}{N}\underset{jl}{}e^{ik(x_jx_l)}\rho _{jl},$$
(19)
which satisfies the relation
$$n_k(N_b)=n_k(NN_b)+\left(1\frac{NN_b}{N/2}\right).$$
(20)
In contrast to periodic systems where the natural orbitals \[Eq. (1)\] are momentum states forrester03 ; rigol04\_1 , this is not the case in a box. (The system is not translationally invariant.) In Fig. 1 we show the lowest-natural-orbital wave function in a box at different temperatures. We have normalized it as
$$\phi ^0=R^{1/2}\varphi ^0,R=\left(N_bN\right)^{1/2},$$
(21)
so that $`\phi ^0`$ vs $`x/N`$ is independent of the system size when the density $`\rho =N_b/N`$ is kept constant.
Figure 1 shows that at finite temperatures the weight of the lowest natural orbital increases in the center of the system, departing from the constant value it would have in the periodic case ($`k=0`$ state). Still, we find that qualitatively (and quantitatively) the natural orbital occupations behave very similarly to the occupations of the momentum states so that for the box we will restrict our analysis to $`n_k`$. The natural orbitals will be relevant to the discussion in the harmonic trap where their behavior can be qualitatively different to the one of $`n_k`$.
In Figs. 2(a)โ2(c) we show the HCB momentum distribution function for half-filled systems with $`N=1000`$ and different temperatures. We have plotted as dashed lines the ground-state results for comparison. The effects of small but finite temperatures are dramatic. This can be better seen in Figs. 2(a) and 2(b) where the energies of the finite-temperature systems are almost identical to the ones of the ground state. For $`k_BT=0.01t`$, the relative energy difference $`(\delta E=[E(T)E(0)]/|E(0)|)`$ between the finite-temperature system \[$`E(T)`$\] and the ground state \[$`E(0)`$\] is $`\delta E0.01\%`$. In Fig. 2(a) one can see that the $`k=0`$ momentum peak is already around 2/3 of the one at zero temperature. For the case in Fig. 2(b), $`\delta E0.4\%`$ and the peak at $`n_{k=0}`$ has already reduced almost 5 times. At $`k_BT=0.5t`$ in Fig.(a), the zero momentum peak has practically disappeared.
As opposed to the HCB momentum distribution function, we have plotted in Figs. 2(d)โ2(f) the momentum distribution function of the equivalent noninteracting fermions. These figures not only show the differences between the shape of the momentum distributions in both cases, but also the fact that they are affected very differently by the temperature. In the fermionic case it is well known that the changes on $`n_k`$ occur only around the Fermi surface and are of order $`k_BT`$, so that in Fig. 2(d) one cannot notice the differences between the finite- and zero-temperature cases. In Fig. 2(e) they are very small, and only when $`k_BT`$ becomes of the order of $`t`$ \[Fig. 2(f)\] can one see a large deviation of the finite-temperature $`n_k`$ with respect to the one in the ground state.
The zero-temperature peaks in the HCB $`n_k`$ \[Figs. 2(a)โFigs. 2(c)\] reflect the presence of quasi-long-range one-particle correlations mccoy68 ; vaidya78 ; mccoy83 ; rigol04\_1 ; i.e., there is a power-law decay $`\rho _{ij}|x_ix_j|^{1/2}`$. In these 1D systems any finite temperature generates an exponential decay of $`\rho _{ij}`$, which destroys the quasi-long-range correlations present in the ground state. This exponential decay is the one producing dramatic effects in $`n_k`$.
In Fig. 3 we show the decay of one-particle correlations for the same systems of Fig. 2. At very low temperatures ($`k_BT=0.01t`$) the one-particle density matrix follows the ground-state result over a certain distance, which reduces with increasing the temperature, approximately up to the point where the exponential decay sets in. Our results in Fig. 3 can be compared with the ones obtained by other means for $`S_i^xS_j^x`$ in the 1D spin-1/2 isotropic XY model tonegawa81 , to which HCBโs can be mapped. Apart from a (1/2) normalization factor the results agree.
The quantity of relevance to characterize the finite-temperature exponential decay of the one-particle density matrix $`\rho _{ij}e^{|x_ix_j|/\xi }`$ (Fig. 3) is the correlation length $`\xi `$. This quantity is of experimental importance since for $`\xi N`$ the HCB gas (essentially) exhibits at finite temperatures properties of the ground state. At low temperatures, $`k_BT<t`$, the correlation length decreases as $`\xi 1/T`$ with increasing temperature. This is shown in Fig. 4.
A way of seeing the effects that a finite-temperature correlation length produces in these bosonic systems is to study how the occupation of the zero-momentum state scales with the number of particles (or the system size) when the density is kept constant. Results for $`n_{k=0}`$ vs $`N_b`$ are presented in Fig. 5. There we have plotted results for as many temperatures as in Fig. 3 so that one can see at what system size the finite-temperature results depart from the ones of the ground state. Since the system size is twice the number of particles (they are at half filling), one can then notice, with the help of Fig. 4, that the mentioned departure indeed occurs for system sizes larger than the correlation length. For example, a half-filled box with 20 HCBโs (a filling similar to the one achieved experimentally in Ref. paredes04 ) would have a momentum distribution function very similar to the one in the ground state up to a temperature $`k_BT=0.05t`$. At zero temperatures $`n_{k=0}`$ scales proportionally to $`\sqrt{N_b}`$; i.e, it diverges when $`N_b\mathrm{}`$, reflecting the power-law decay of one-particle correlations shown in Fig. 3 rigol04\_1 .
The one-particle correlation length not only depends strongly on the temperature, but also on the density in the system. (In Fig. 4 we have only shown results for the half-filled case.) The dependence of the correlation length on the density, for two values of the temperature, is depicted in Fig. 6. Notice that both curves are symmetric with respect to $`\rho =0.5`$ due to particle-hole symmetry.
The strong dependence of the correlation length on the density represents a difficulty for defining $`\xi `$ in inhomogeneous systems, like the ones achieved experimentally where HCBโs are trapped in harmonic confining potentials. (In a box the density is not exactly constant, away from half and integer fillings, due to Friedel oscillations, but they reduce with increasing system size.) An alternative definition to the correlation length $`\xi `$ may be given as the second moment of the one-particle density matrix
$$\stackrel{~}{\xi }=\sqrt{\frac{1}{2}\frac{_{ij}\left(x_ix_j\right)^2\rho _{ij}}{_{ij}\rho _{ij}}}.$$
(22)
In Fig. 4 we have plotted $`\stackrel{~}{\xi }`$ along with $`\xi `$. When $`1<\xi N`$ both $`\stackrel{~}{\xi }`$ and $`\xi `$ are very similar. For the lowest temperatures, in Fig. 4, we considered systems with 1000 lattice sites, which are not much larger than the correlation length. This is the origin of the differences between $`\stackrel{~}{\xi }`$ and $`\xi `$ observed for large values of $`\xi `$. At high temperatures ($`k_BT>t`$) the value of $`\stackrel{~}{\xi }`$ is completely dominated by the very-short-distance sector of the one-particle density matrix, so that $`\stackrel{~}{\xi }`$ and $`\xi `$ are expected to be very different. At intermediate temperatures one can use $`\stackrel{~}{\xi }`$ as a good estimate of $`\xi `$.
In Fig. 6 we have also plotted $`\stackrel{~}{\xi }`$ along with $`\xi `$ so that one can realize how the inclusion of the short-range part of the one-particle density matrix in $`\stackrel{~}{\xi }`$ produces different effects for low densities, where $`\stackrel{~}{\xi }>\xi `$, and high densities, where $`\stackrel{~}{\xi }<\xi `$. Still the overall behavior of $`\stackrel{~}{\xi }`$ is similar to the one of $`\xi `$. In the next section we will rely on $`\stackrel{~}{\xi }`$ for estimating the correlation length in harmonic traps and also for comparing it to the one in the box.
## IV Hard-core bosons in harmonic traps
We study in this section HCBโs trapped in harmonic potentials. The addition of a confining potential generates a position-dependent density profile where, at zero temperature, superfluid and Mott-insulating regions can coexist. In the next two subsections we analyze the effects of the temperature on density and momentum profiles of systems in which the ground state is (i) superfluid (Sec. IV.1) and (ii) a coexistence of superfluid and Mott-insulating phases (Sec. IV.2). In Sec. IV.3 we address more general questions like the behavior of one-particle correlations and scaling properties at finite temperatures.
In harmonic traps we normalize $`n_k`$ using a length scale set by the combination lattice-confining potential,
$$\zeta =\left(V_2/t\right)^{1/2},$$
(23)
so that
$$n_k=\frac{a}{\zeta }\underset{jl=1}{\overset{N}{}}e^{ik(jl)}\rho _{jl}.$$
(24)
In addition, instead of the density $`\rho =N_b/N`$, relevant to the periodic or open case, we consider the characteristic density rigol04\_1 ; rigol03\_3
$$\stackrel{~}{\rho }=N_ba/\zeta .$$
(25)
As shown in Ref. rigol03\_3 up to $`\stackrel{~}{\rho }2.6`$โ2.7 there is no Mott insulator in the trap. For larger values of $`\stackrel{~}{\rho }`$ a Mott-insulating phase appears in the middle of the system.
### IV.1 Superfluid case at $`T=0`$
In Fig. 7 we show density and momentum profiles in a trap with 200 HCBโs ($`\stackrel{~}{\rho }=2`$) for different temperatures (solid line) and compared to the ground-state case (dashed line). As for the fermionic $`n_k`$ in Figs. 2(d)โ2(f), the changes of the density profiles with the temperature in Figs. 7(a)โ7(c) are the ones expected for fermions. \[HCBโs and fermions exhibit identical density profiles, Eq. (II).\] For temperatures much smaller than the Fermi energy, which is of the order of $`t`$ for these systems, the density profiles almost do not change. The same occurs with the total energy $`E`$ of the trapped cloud, as seen from their values reported in the caption of Fig. 7. On the other hand, the behavior of $`n_k`$, related to off-diagonal one-particle correlations, is very different to the one of the density. At $`k_BT=0.01t`$ \[Figs. 7(a) and 7(d)\], when the energy of the system has changed by $`0.03\%`$ with respect to the ground-state energy, changes can be already noticed in $`n_k`$ around $`k=0`$. For $`k_BT=0.1t`$ \[Figs. 7(b) and 7(e)\], the energy is $`3\%`$ larger than in the ground state and the peak in $`n_{k=0}`$ is less than one-third of its value at $`T=0`$. For larger temperatures, like $`k_BT=0.5t`$ in Figs. 7(c) and 7(f), almost no peak can be seen in $`n_{k=0}`$ as compared with the one in the ground state. This is similar to the results obtained for the box in the previous section, with a difference being that in the box the density distribution is not affected by the temperature.
Other quantities of relevance to the harmonically trapped case are the natural orbital occupations \[Figs. 8(a)โ8(c)\] and the wave function of the lowest natural orbital \[Figs. 8(d)โ8(f)\]. In Figs. 8(d)โ8(f) we normalize the natural orbital wave function following Ref. rigol04\_3
$$\phi ^0=R^{1/2}\varphi ^0,R=\left(N_b\zeta /a\right)^{1/2}.$$
(26)
Like $`n_k`$, the natural orbital occupations exhibit a very strong dependence on the temperature, which can be understood since they are also related to the off-diagonal one-particle correlations. More interesting, and qualitatively different to the case in the box, is the behavior displayed by the lowest-natural-orbital wave function in Fig. 8. With increasing temperature, for large fillings, the weight of the lowest natural orbital in the middle of the trap decreases, and for $`k_BT=0.5t`$ \[Fig. 8(f)\] it is exactly zero. This behavior of the wave function is accompanied by the appearance of a degeneracy in the occupation of the lowest natural orbitals. These two effects are very similar to the ones generated by the increase of the filling in the ground state of the system and the formation of a Mott insulator in the middle of the trap rigol04\_1 . However, as seen in Figs. 7(a)โ7(c) no Mott insulator is created by an increase of the temperature.
In order to understand the above effect it is important to realize that the spectrum of noninteracting particles in a combination lattice-harmonic potential rigol03\_3 ; hooley04 ; ruuska04 ; rey05 is very different to the one of the harmonic oscillator in the continuum. (In the latter case one can intuitively realize that the maximum weight of a condensate, or of the largest eigenvalue of the one-particle density matrix, occurs in the middle of the trap.) In a lattice with a superposed harmonic oscillator, the eigenvalues of the noninteracting Hamiltonian reduce their weight in the middle of the system when their energy increases. For energies larger of $`2t`$ \[for the Hamiltonian in Eq. (5)\] BW , the eigenstates of the Hamiltonian start to be localized at the sides of the trap; i.e., they have zero weight in the center of the system rigol03\_3 ; hooley04 ; ruuska04 ; rey05 . At zero temperatures a Mott-insulating domain in the center of the trap signals that these states are populated rigol03\_3 . At finite temperatures the occupation of localized states occurs even when there is no Mott insulator in the system, which explains why the lowest natural orbital can exhibit a behavior like the one seen in Figs. 8(e) and 8(f) in the absence of the insulating core.
Before analyzing the finite-temperature one-particle correlations in the confined system, which explain the previously observed effects in $`n_k`$ and the natural orbital occupations, we present in what follows an example of the consequences of the temperature in a system that in its ground state exhibits a coexistence of superfluid and Mott-insulating phases.
### IV.2 Mott insulator is present at $`T=0`$
In Fig. 9 we show density and momentum profiles of a system with 300 HCBโs ($`\stackrel{~}{\rho }=3`$) for different temperatures (solid line) and compared to the ground-state case (dashed line). As for the superfluid case discussed in the previous subsection, density profiles are almost not modified for $`k_BTt`$. Increasing the temperature one can see in Fig. 9(b) that as $`k_BT`$ approaches $`t`$ the Mott insulating ($`n=1`$) plateau in the middle of the trap disappears. The effects of the temperature in $`n_k`$ are also similar to the ones in the case with no Mott insulator. $`n_k`$ strongly depends on the temperature in the system.
It is worth noticing in Figs. 9(c) and 9(d) that even in the presence of a Mott-insulating phase, at zero temperature, $`n_k`$ exhibits a sharp peak at $`k=0`$ due to the superfluid phases at the sides wessel04 ; rigol04\_1 . The effects of the Mott insulator in $`n_k`$ are reflected by a large population of $`k`$ states around $`ka=\pm \pi `$ and an increase of the full width at half maximum of the $`k=0`$ peak. These are characteristics of the system that remain at finite but very low temperatures. Increasing the temperature \[Fig. 9(d)\] the peak at $`k=0`$ disappears. On the other hand, the high-momentum tails remain almost unmodified with respect to the ground-state case as they reflect the properties of short-distance correlations, related to the density profiles, which are much less sensitive to temperature effects.
At zero temperature, the natural orbital occupations exhibit a clear signature of the presence of the Mott-insulating core in the trap. A plateau with $`\lambda _\eta =1`$ is present, reflecting the existence of single occupied states. In this case the lowest natural orbital is degenerate due to the splitting of the system by the Mott-insulating core. Two identical quasicondensates can be observed at the sides of the Mott core \[Figs. 10(c) and 10(d)\]. The increase of temperature reduces the occupation of the lowest natural orbital \[they become more localized, Figs. 10(c) and 10(d)\], but does not destroy their degeneracy. This degeneracy in absence of a Mott-insulating state is, as explained in the previous subsection, an effect that only appears at finite temperatures due to the population of localized states at the sides of the trap. Finally, one should notice that the plateau with $`\lambda _\eta =1`$ disappears in Fig. 10(b) along with the disappearance of the Mott plateau in Fig. 9(b).
### IV.3 Correlation functions and scalings
In Fig. 11 we show the behavior of one-particle correlations with increasing temperature for the two cases analyzed in the previous subsections. Correlations ($`\rho _{ij}`$) are measured with respect to a fixed point $`x_j`$, while $`x_i`$ is changed all over the system. In Fig. 11(a) the correlations are measured with respect to the middle of the trap and in Fig. 11(b) with respect to two points at the sides of the Mott-insulating core present at $`T=0`$.
In the ground state, the one-particle density matrix decays as a power law $`\rho _{ij}|x_ix_j|^{1/2}`$ for $`0<n_i,n_j<1`$ rigol04\_1 . The introduction of a small temperature $`k_BT=0.01t`$ can be already noticed in Fig. 11(a) as a faster decay of correlations at long distances. At temperatures larger than $`k_BT=0.1t`$, for the system sizes of the figure, the one-particle density matrix decays exponentially to $`10^8`$ before reaching the borders of the trap. Due to the space varying density one can notice that in contrast to the box, in a harmonic trap one cannot see a single correlation length, which would mean a straight line in all the semilogarithmic plots of the figure. Still one can calculate the second moment of the one-particle density matrix $`\stackrel{~}{\xi }`$ \[Eq. (22)\] as a sort of an averaged correlation length.
We present in Table 1 results for $`\stackrel{~}{\xi }`$ in harmonic traps for two temperatures and six values of $`\stackrel{~}{\rho }`$. To the right we show results obtained in boxes with densities chosen to be identical to the ones at the center of the harmonically trapped cloud. One can see that the results in both cases are similar far from the region where the Mott insulator sets in the middle of the trap ($`\stackrel{~}{\rho }=`$0.5โ2.0 in Table 1), so that one can estimate $`\stackrel{~}{\xi }`$ in harmonic traps using results from a box. This is in agreement with recent results reported for other finite-temperature correlation lengths in trapped bosonic systems with no lattice kheruntsyan05 . The reason for the agreement between $`\stackrel{~}{\xi }`$ in the trap and in the box is that in the first case $`\stackrel{~}{\xi }`$ is dominated by the contributions of the middle of the system, where the density has its maximum value and it is โrelatively uniform.โ As one approaches $`n=1`$ in the middle of the trap, or $`\stackrel{~}{\rho }=2.6`$, the argument above fails ($`\stackrel{~}{\rho }=2.5,3.0`$ in Table 1) since the correlation length in the center of the cloud approaches zero (see $`\stackrel{~}{\xi }`$ vs $`\rho `$ in Fig. 6 when $`\rho 1`$) and regions with smaller densities start to dominate the value of $`\stackrel{~}{\xi }`$. For those cases an exact calculation of $`\stackrel{~}{\xi }`$, given the density profile, is required.
The exponential decay of the one-particle density matrix implies that when the size of the system is larger than the averaged correlation length, the momentum distribution function and the occupation of the lowest natural orbital stop changing with increasing system size. The size at which this occurs depends on the temperature, as the averaged correlation length decreases with increasing the temperature, and also depends on the characteristic density. In Fig. 12 we show how $`n_{k=0}`$ and $`\lambda _0`$ scale at four different temperatures and starting from small system sizes (close to the ones achieved experimentally paredes04 ). At $`T=0`$ the increase of both quantities is $`\sqrt{N_b}`$, reflecting quasi-long-range correlations present in $`\rho _{ij}`$ rigol04\_1 . For $`k_BT=0.01t`$, the departure from the zero-temperature values occurs when the trap has $`100`$ HCBโs. For $`k_BT=0.05t`$ it occurs around $`N_b=20`$, and for $`k_BT=0.1t`$ even the smallest system with 10 HCBโs is different to the ground state.
To conclude this section we show in Fig. 13 a comparison between density and momentum profiles for 100 and 400 HCBโs at $`k_BT=0.1t`$ and $`\stackrel{~}{\rho }=2`$. While even at zero temperature the density profiles do not differ rigol04\_1 ; rigol03\_3 , the ground-state peak $`n_{k=0}`$ would have been 2 times larger for 400 HCBโs than for 100 HCBโs. At $`k_BT=0.1t`$ both momentum distribution functions are almost indistinguishable.
## V Grand-canonical vs canonical ensemble
In the previous sections we have discussed effects of the temperature on trapped HCBโs in 1D. The starting point for our calculations was the grand-canonical ensemble, in which the system is assumed to be in thermal equilibrium with a large reservoir with temperature $`T`$ and chemical potential $`\mu `$. The chemical potential was then chosen to obtain the desired average number of particles in the trap. In this section we analyze the changes introduced by the grand-canonical fluctuations of the particle number with respect to a fixed-$`N_b`$ canonical description, which may be more relevant to describe trapped ultracold quantum gases where no particle reservoir is available.
In the thermodynamic limit both descriptions are known to provide the same predictions huang87 . On the other hand, for noninteracting bosonic systems with a mesoscopic number of particles, it has been shown that the differences between the grand-canonical and canonical condensate fractions can be as large as $`10\%`$ (for $`N_b=100`$) close to the BEC transition point and decreasing logarithmically with increasing number of particles. In the present work we have been dealing with the opposite caseโi.e., infinite repulsion. Interactions are known to suppress fluctuations of the number of particles in the grand-canonical ensemble huang87 , but since recent experiments with HCBโs on optical lattices achieved only up to 20 HCBโs in around 50 lattice sites, it is useful to present an estimate of the difference between both ensembles for such small systems.
In order to obtain the canonical one-particle density matrix we use the ground-state approach of Ref. rigol04\_1 . We calculate the Greenโs function of all states $`|\mathrm{\Psi }_{HCB}^n`$ with $`N_b`$ bosons in $`N`$ lattice sitesโi.e., of $`N_s=N!/(NN_b)!N_b!`$ states. The canonical Greenโs function at temperature $`T`$ is obtained as the sum
$$G_{ij}^C=\frac{1}{Z^C}\underset{n=1}{\overset{N_s}{}}e^{E_n/k_BT}\mathrm{\Psi }_{HCB}^n|b_ib_j^{}|\mathrm{\Psi }_{HCB}^n,$$
(27)
where $`e^{E_n/k_BT}`$ ($`E_n`$ is the energy of state $`|\mathrm{\Psi }_{HCB}^n`$) is the Boltzmann factor and $`Z^C`$ the canonical partition function:
$$Z^C=\underset{n=1}{\overset{N_s}{}}e^{E_n/k_BT}.$$
(28)
The canonical one-particle density matrix is then
$`\rho _{ij}^C`$ $`=`$ $`{\displaystyle \frac{1}{Z^C}}{\displaystyle \underset{n=1}{\overset{N_s}{}}}e^{E_n/k_BT}\mathrm{\Psi }_{HCB}^n|b_i^{}b_j|\mathrm{\Psi }_{HCB}^n`$ (29)
$`=`$ $`G_{ij}^C+\delta _{ij}\left(12G_{ii}^C\right).`$
In Fig. 14 we show results obtained for the grand-canonical ($`E`$) and canonical ($`E^C`$) energies of 10 HCBโs in a box with 50 lattice sites as a function of the temperature. At the scale of the figure they are indistinguishable. More information can be obtained in the inset (a) where we plot as thin lines the energy difference between both ensembles as a function of the number of particles in boxes with densities $`N_b/N=0.2`$. As seen in this inset even for such small systems the difference almost does not change with $`N_b`$, and it is always smaller than the energy unit $`t`$. Considering that the modulus of the energy increases linearly with the system size, the relative difference between both ensembles decreases $`\delta E(EE^C)/|E^C|1/N_b`$ \[thick lines in inset (a)\]. For $`N_b=10`$ and $`N=50`$ one can see that $`\delta E`$ is below $`1\%`$ for temperatures up to $`k_BT=0.5t`$. As the temperature increases beyond $`k_BT=t`$ the differences between $`E`$ and $`E^C`$ start to decrease, which together with the decrease of the modulus of $`E^C`$ shown in Fig. 14 produces a saturation of $`\delta E`$ at around $`2\%`$ for $`k_BT>10t`$. Then for $`N_b=10`$ and $`N=50`$ the maximum $`\delta E`$ is just a $`2\%`$ of the energy. We have also studied other densities keeping $`N=50`$, and the results obtained for the maximum $`\delta E`$ were exactly the same $`2\%`$.
While Kinoshita et al. kinoshita04 used the energy of the system to confirm the achievement of the hard-core limit, Paredes et al. paredes04 considered the momentum distribution function. In the inset (b) of Fig. 14 we show the relative difference \[$`\delta n_k(_k|n_kn_k^C|)/(_kn_k^C)`$\] between the grand-canonical $`n_k`$ and canonical $`n_k^C`$ calculation of the momentum distribution function. The relative differences for $`n_k`$ although larger than the corresponding ones for the energy are still small and also reduce $`1/N_b`$ with increasing the system size. For $`N_b=10`$ and $`N=50`$ they are smaller than $`1\%`$ up to $`k_BT=0.2t`$.
It is also useful to calculate the differences between the grand-canonical and canonical ensemble for the equivalent noninteracting fermions. This may be relevant for systems like the ones recently achieved experimentally by Kรถhl et al. kohl05 . For noninteracting fermions, the energy differences between both ensembles are identical to the ones of the HCBโs due to the mapping, Eqs. (2)โ(5), so that as shown in Fig. 14 and its inset (a) they are small. For the fermionic momentum distribution function the HCB results do not apply. We have also calculated the fermionic relative difference \[$`\delta n_k^f(_k|n_k^fn_k^{f,C}|)/(_kn_k^{f,C})`$\] between the grand-canonical $`n_k^f`$ and canonical $`n_k^{f,C}`$ calculation of $`n_k`$. They are larger than for the ones of the HCBโs as shown in the inset (b) of Fig. 14. However, they are still small for the experimentally accessible system sizes. For $`N_f=10`$ and $`N=50`$ they are smaller than $`3\%`$ for $`k_BT=0.2t`$. Apart from an even-odd effect that decreases with increasing the temperature, $`\delta n_k^f`$ also decreases $`1/N_b`$ with increasing the system size.
The introduction of a harmonic trap does not (qualitatively) change the results obtained in a box. In Fig. 15 we show the grand-canonical and canonical results of the energy in a harmonic trap with 10 particles and $`\stackrel{~}{\rho }=2`$ as a function of the temperature. Contrary to the box, in a harmonic trap the energy is not bounded from above for very large temperatures. This is because the HCB cloud can increase its size and consequently its potential energy. The energy differences between both ensembles, when changing the number of particles keeping $`\stackrel{~}{\rho }=2`$ constant, are shown in the inset of Fig. 15. As for the box they are almost independent of $`N_b`$, and smaller than $`t`$. The results for the relative differences between the momentum distribution functions for HCBโs and noninteracting fermions in both ensembles are also shown in the inset. Their behavior is very similar to the one of the box in inset (b) of Fig. 14.
As mentioned before, Herzog and Olshanii herzog97 discussed the grand canonical and canonical differences between the condensate fraction for noninteracting bosons in harmonic traps. At finite repulsive interactions, in 1D, there is no BEC even at zero temperature. Still, for the HCBโs we have calculated the differences between the largest eigenvalue of the one-particle density matrix (equivalent to the condensate occupation for BEC leggett01 ) in the grand-canonical and canonical ensembles. As for $`\delta n_k`$ in the inset of Fig. 15, we find that the difference between them decreases $`1/N_b`$ with increasing number of particles in the system. This clearly contrasts with the $`1/\mathrm{ln}(N_b)`$ obtained for the noninteracting case herzog97 .
We conclude by explicitly showing in Fig. 16 the density profiles and momentum distribution functions of 10 HCBโs in a harmonic trap with $`\stackrel{~}{\rho }=2`$ at $`k_BT=0.5t`$ as obtained from the grand-canonical and canonical descriptions. They are basically indistinguishable. Then, even for the small system sizes achieved experimentally, one can rely on the grand-canonical description for strongly correlated HCBโs for the physical quantities described here. The same conclusion applies to noninteracting fermions.
## VI Conclusions
We have presented an exact study of finite-temperature properties of HCBโs confined on 1D lattices. In order to solve this problem we have used the Jordan-Wigner transformation to map HCBโs into noninteracting fermions. After the mapping, properties of Slater determinants allowed us to obtain an exact expression for the HCB one-particle density matrix in terms of determinants of $`N\times N`$ matrices, which are evaluated numerically. Our approach represents an alternative for finite systems to previous works that considered the thermodynamic limit lieb61 ; mccoy68 ; vaidya78 ; mccoy83 ; tonegawa81 for periodic and open chains and in which Toeplitz determinants were involved.
We have shown that the effects of small finite temperatures are very important when dealing with quantities related to off-diagonal one-particle correlations like the momentum distribution function and the natural orbitals. These finite-temperature effects depend strongly on the system size. On the other hand, observables related to diagonal one-particle correlations (identical for fermions), like density profiles, are much less affected at low temperatures. Explicit results for the behavior of all these quantities versus temperature were given for system sizes that range from the ones recently achieved experimentally up to 20 times larger.
Finally, we have compared grand-canonical and canonical results for energies and momentum distribution functions of HCBโs and noninteracting fermions for small systems, like the ones achieved experimentally. In spite of the mesoscopic number of particles we have shown that for these system sizes the effects of the grand-canonical fluctuations of the particle number are very small and one can rely on a grand-canonical approach.
Although all our calculations are exact for infinite on-site $`U`$ repulsion, for very strong but finite $`U`$ our conclusions are still valid. In this case $`1/U`$ acts like a perturbation to the noninteracting spinless fermion Hamiltonian cazalilla03 . Rey et al. rey05 have recently presented results obtained by exact diagonalization that support the above conclusion rey05 . A connection to experimentally relevant parameters can be also found in Ref. rey05 .
###### Acknowledgements.
We are grateful to G. G. Batrouni, A. Muramatsu, M. Olshanii, R. T. Scalettar, and R. R. P. Singh for stimulating discussions and comments on the manuscript and to M. Arikawa for pointing out several references. This work was supported by Grant Nos. NSF-DMR-0312261, NSF-DMR-0240918, and NSF-ITR-0313390.
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# Discovery of a redshifted X-ray emission line in the symbiotic neutron star binary 4U 1700+24
## 1 Introduction
Based on its optical identification with the M2 III giant HD154791 (Garcia et al. 1983; Masetti et al. 2002; also known as V934 Her), the X-ray source 4U 1700$`+`$24 is classified as both a Low Mass X-ray Binary (Liu et al. 2001) and a symbiotic-like binary (Garcia et al. 1983). The lack of signatures of binarity and the initially marginal positional coincidence (Morgan & Garcia 2001) originally made this identification uncertain. More recently, the detection of radial velocity variations from the proposed optical counterpart with the same period as X-ray brightness variations (Galloway et al. 2002) has added support to the identification.
In X-rays, 4U 1700$`+`$24 has long periods in which it is faint, and occasional episodes in which the flux increases by more than an order of magnitude. Assuming a distance of 420 pc (Masetti et al. 2002), the 2โ10 keV luminosity of 4U 1700$`+`$24 varies from 2$`\times `$10<sup>32</sup>, which is much higher than expected from an isolated M giant of this spectral type, to 10<sup>34</sup> erg s<sup>-1</sup>. Masetti et al. (2002) reported the analysis of six observations from five different X-ray satellites, including an observation with the Rossi X-Ray Timing Explorer (RXTE; Jahoda et al. 1996) taken during a $`100`$-day outburst in November 1997. Although no coherent pulsations or QPOs were detected, their analysis shows that the X-ray flux is also variable on short time scales ($``$ seconds).
In July 2002, the RXTE All-Sky Monitor light curve of 4U 1700+24 indicated that the source was undergoing a new outburst, and a Target of Opportunity observation was performed by the XMM-Newton satellite. In the following sections we discuss the analysis of the spectroscopic X-ray data from this observation.
## 2 Observations and data reduction
4U 1700$`+`$24 was observed with the XMM-Newton satellite (Jansen et al. 2001) on 2002 August 11. The XMM-Newton X-ray optics consist of three nested Wolter-I mirror assemblies, illuminating five X-ray detectors which operate simultaneously. CCDs sensitive to photons in the range 0.15-15 keV are situated at the focal point of each the three mirrors, two MOS and one PN types, comprising the European Photon Imaging Camera (EPIC; Strรผder et al. 2001, Turner et al. 2001). Two of the mirrors also illuminate the Reflection Grating Spectrometers (RGS; den Herder et al. 2001) which disperse photons in the range 5โ35 ร
onto a pair of dedicated off-axis CCDs. The PN and MOS cameras have spectral resolutions of $``$80 eV and $``$70 eV (at 1 keV) and point spread functions (PSFs) of $``$6<sup>โฒโฒ</sup> and $``$5<sup>โฒโฒ</sup> for the PN and MOS cameras, respectively. The RGS spectral resolution is 0.04 ร
although it has a significantly smaller effective area ($``$100 cm<sup>2</sup> at 1 keV for each of the two units, compared to $``$1200 cm<sup>2</sup> for the PN and $``$900 cm<sup>2</sup> for the sum of the two MOS detectors) and cannot produce 2-dimensional images.
The exposure time of the 4U 1700$`+`$24 observation (corrected for the deadโtime) was 6 ks for the PN, 7 ks for the MOS, and 8 ks for the RGS. To avoid excessive optical loading on the PN and MOS cameras from the bright optical companion of 4U 1700$`+`$24 (V=7.6), the โthickโ filter was selected. The X-ray brightness of 4U 1700$`+`$24 during this high-state observation necessitated the use of Small Window operating mode, in which only a small portion (about 4$`{}_{}{}^{}\times `$4 for the PN and 2$`{}_{}{}^{}\times `$2 for the MOS) of the CCD is read out in a reduced frame time (6 ms for PN and 300 ms for MOS). The short readout time of this mode helps reduce errors due to photon pileโup, which occurs when two (or more) photons hit nearby pixels during a single CCD exposure, producing an event that is indistinguishable from a single, higher energy photon. The main effects of photon pileโup are an incorrect reconstruction of the event energies and an underestimation of the source count rate.
All the data were processed using the XMM-Newton Science Analysis System (SAS, version 6.1.0)<sup>1</sup><sup>1</sup>1http://xmm.vilspa.esa.es/external/xmm\_sw\_cal/sas\_frame.shtml. The use of Small Window mode only partially mitigated the effects of pileโup<sup>2</sup><sup>2</sup>2The presence of pileโup in a dataset can be investigated by examining the energy distribution of single, double, and multiple events (composed, respectively, of one, two, and greater than two adjacent pixels being above the detection threshold) in the PN and MOS data, and so in our extraction of source photons we excluded the central part of the source PSF, where the incident count rate per pixel is the highest. The EPIC spectra were thus extracted from annular regions centered on the source, with inner and outer radii of 10$`\mathrm{}`$ and 40$`\mathrm{}`$, respectively. With these regions, we were able to completely remove the portion of the image affected by pileโup from the PN data, but not from the MOS data. The choice of an even larger inner radius for MOS would have greatly reduced the number of valid counts. Therefore, the MOS spectra are not discussed in this paper. To generate the PN spectrum, only single and double events were retained, and the resulting spectrum was rebinned both to have at least 30 counts in every energy channel and at most three bins per PN energy resolution element. A background spectrum was extracted from a region far enough from 4U 1700$`+`$24 to avoid contamination by source photons. In the 0.2โ12 keV range, the background contributes less than 0.5% of the total count rate for most of the observation, increasing to 3% during the last 1500 seconds.
The RGS spectra (source and background) were extracted, and response matrices calculated using the standard reduction procedures. Since the source photons are widely spread along the dispersion axis, the RGS data are not affected by pileโup. The RGS source spectra were also rebinned in order to have at least 30 counts in every channel.
## 3 Analysis and results
### 3.1 Spectral Continuum
In order to compare the XMM-Newton spectrum of 4U 1700$`+`$24 with previous observations, we fitted the PN spectrum in the 0.3โ12 keV range with the two-component model used by Masetti et al. (2002): a blackbody plus Comptonization (COMPST; Sunyaev & Titarchuk 1980), with absorption fixed to the Galactic value of $`N_H`$ = 4$`\times `$10<sup>20</sup> cm<sup>-2</sup> (Dickey & Lockman 1990). The resulting fit is unacceptable ($`\chi _{red}^2`$ = 4.05 for 273 degrees of freedom) due to a large excess at low energies. The fit can be improved by adding a component that contributes below $`1`$ keV. We obtain the best results by modeling this soft excess with a broad ($`\sigma `$0.1 keV) Gaussian line centered at $``$0.5 keV. We report the spectral parameters in Table 1. We show the PN spectrum and its residuals with respect to this model in Figure 1. Although even our final fit is only marginally acceptable ($`\chi _{red}^2`$ = 1.19 for 270 d.o.f., corresponding to a null hypothesis probability of 0.018), the addition of a systematic error of 1.5%, which is compatible with the current calibration uncertainties of the PN, reduces $`\chi _{red}^2`$ to 1. No narrow features are detected in the PN spectrum. In particular, the 3$`\sigma `$ upper limit on the equivalent width of a Fe K-$`\alpha `$ fluorescent line at 6.4 keV is $`15`$ eV.
### 3.2 The emission line at 19.2 ร
The high spectral resolution of the RGS allows us to search for narrow features in the X-ray spectrum of 4U 1700$`+`$24 at low energies. In Figure 2, we show the first-order RGS spectrum of 4U 1700$`+`$24 and residuals with respect to the best-fit model derived from analysis of the PN data. In addition to some deviations from the continuum model due to problems in the cross-calibration between the PN and RGS instruments (Kirsch et al. 2004), an emission line is apparent at about 19 ร
(Fig. 2). The addition of a Gaussian component in the model accounts for the residuals in this region of the spectrum (see bottom panel of figure 2) and improves the $`\chi ^2`$ from 1.65 (for 765 d.o.f.) to 1.54 (for 762 d.o.f.). The emission-line parameters are: $`\lambda `$ = 19.19$`{}_{0.09}{}^{}{}_{}{}^{+0.05}`$ ร
, $`\sigma `$ = $`3.9_{1.3}^{+2.7}`$ eV, line flux equals $`(4.9_{1.6}^{+1.9})\times 10^4`$ photons cm<sup>-2</sup> s<sup>-1</sup>, and equivalent width equals 9.7$`{}_{3.1}{}^{}{}_{}{}^{+3.8}`$ eV (all the quoted errors are 3$`\sigma `$ errors).
In the Xโray atomic lines catalogue ATOMDB 1.3.1<sup>3</sup><sup>3</sup>3http://asc.harvard.edu/atomdb/, only very low emissivity transitions are found within 3$`\sigma `$ of the line position; most are satellite lines of H-like and He-like oxygen, with a maximum emissivity of 8.4$`\times 10^{18}`$ photons cm<sup>3</sup> s<sup>-1</sup>. The Ly-$`ฯต`$ line of N VII (which is not included in the ATOMDB database) has a wavelength compatible with the observed emission line ($`\lambda `$ = 19.118 ร
, Verner et al. 1996a), but the probability for this transition is very low. The Ly-$`\alpha `$ transition of H-like oxygen (O VIII) has an emissivity more than 2 orders of magnitude higher than that of the lines mentioned above (which makes it typically one of the strongest lines found in cosmic Xโray sources) and a rest-frame wavelength of 18.97 ร
, slightly smaller than the value we find for the emission line shown in Figure 2. It is therefore possible that the X-ray emission line in 4U 1700$`+`$24 is O VIII at redshift $`z=0.012_{0.004}^{+0.002}`$.
The accuracy of the RGS wavelength scale is better than 10 mร
(Pollock 2004), with an additional 2.3 mร
of systematic error for each arcsecond error in the pointing or source coordinates. For the 4U 1700$`+`$24 observation, the additional error is smaller than 10 mร
, since the EPIC positional accuracy is better than 4<sup>โฒโฒ</sup> and the RGS data were processed using source coordinates derived from the EPIC images. Therefore, we exclude the possibility that the wavelength shift from the O VIII rest-frame position is due to an incorrect wavelength scale in the RGS instrument.
Other smaller structures are also present in the residuals. In particular, the data differ substantially from the model around the instrumental oxygen edge at $`23.5`$ ร
, where the RGS effective area calibration is complex. We can improve the RGS spectral fit slightly ( $`\chi _{red}^2`$=1.51 for 761 d.o.f.) by including an overabundance of neutral oxygen in the photoelectric absorption model (with cross section from Verner at al. 1996b). Keeping the hydrogen column density fixed at the value of 4$`\times `$10<sup>20</sup> cm<sup>-2</sup>, we obtain a best-fit value for the oxygen abundance of 2.0$`{}_{0.5}{}^{}{}_{}{}^{+0.3}`$ times the solar value (Anders & Grevesse 1989). On the other hand, two instrumental absorption features at 23.05 and 23.35 ร
contribute $`2\times `$10<sup>17</sup> cm<sup>-2</sup> to the oxygen column density (de Vries et al. 2003) and could also explain most of the detected overabundance of oxygen.
If the emission line at 19.19 ร
is redshifted O VIII, we might expect spectral features from other ionization states of oxygen. The signature of He-like oxygen is a triplet of X-ray lines at 21.6, 21.8, and 22.1 ร
. Some residuals are in fact present in the RGS spectrum around 21.8 ร
, which corresponds to either the rest-wavelength of the intercombination line or a resonance line at $`z`$=0.01. Fitting the residuals at 21.8 ร
with a Gaussian of the same width as the O VIII line candidate ($`\sigma `$ = 4 eV) gives a 3$`\sigma `$ upper limit for the equivalent width of 5 eV. Fully ionized oxygen should produce an O VIII radiative recombination continuum at around 14 ร
(14.2 ร
for plasma at rest and 14.4 ร
at $`z=0.01`$). The RGS spectrum shows some residuals at this location. Including a recombination emission edge at 14.4 ร
with a 50 eV width in the fit gives a normalization of (6.1$`\pm `$3.1)$`\times `$10<sup>-4</sup> photons cm<sup>-2</sup> s<sup>-1</sup>. The addition of this component is statistically significant, although it is quite broad and could also result from uncertainties in the spectral continuum model. The RGS spectra were also rebinned using different criteria, but no other spectral features could be identified.
Due to the poorer energy resolution of the PN camera, no line at 19.19 ร
is significantly detected in the PN data. However, if we fix the energy and the width of the line to the values derived from the RGS data, the upper limit to the line equivalent width is 12 eV, which is consistent with the RGS results.
## 4 Discussion
The high sensitivity of XMM-Newton at low energies and the high spectral resolution of the RGS instrument have revealed new features in the X-ray spectrum of the unusual interacting binary 4U 1700$`+`$24 . Although the PN spectrum above 1 keV is consistent with the rather variable spectra seen in previous observations (Masetti et al. 2002), we clearly detect both a soft excess and an emission line at 19.2 ร
.
The best fit to the soft excess was found by adding a broad Gaussian line component to the high-energy continuum model (blackbody + Comptonization). But the interpretation of this feature as an emission line is problematic. No strong emission lines are expected at the wavelength of the soft excess for either $`z=`$0.01 or $`z=`$0. In addition, both the intensity and broadening of this line would be much larger than for the line at 19.2 ร
.
As mentioned above, we have also discovered an emission line at 19.2 ร
. Since only very low emissivity lines are consistent with the observed wavelength of $`\lambda `$ = 19.19$`{}_{0.09}{}^{}{}_{}{}^{+0.05}`$ ร
, the line may be O VIII at redshift $``$0.01. For most high-emissivity lines close to this wavelength, a much larger redshift/blueshift would be required. For example, among the other strong lines expected in the spectrum of an accreting X-ray binary, the O VII and Ne IX triplets are the nearest candidates for the line identification, at longer and shorter wavelengths, respectively. For O VII, the line would be blue-shifted by a factor $`>0.1`$, which, in case of Doppler shift, requires that the emitting plasma is moving towards us at more than 10% of the speed of light. On 2002 July 29 radio observations were performed to look for a possible jet, but no source was detected with an upper limit of 1.0$`\pm 0.7`$ mJy at 15 GHz (G. Pooley 2002, private communication).
On the other hand, the identification of the line at 19.2 ร
with the Ne IX triplet would imply redshifts of 0.40, 0.41, and 0.42 for the forbidden, intercombination and resonance lines, respectively. These values could be interpreted as gravitational redshift from close to the surface of a neutron star, as they are consistent with most of the equations of state for neutron stars composed of normal nuclear matter, and they are just slightly larger than the redshift $`z`$=0.35 found by Cottam et. al. (2002) in the spectral analysis of X-ray bursts from EXO 0748โ676. However, due to the lack of identifications of other spectral features at $`z`$ 0.4 or indication for overabundance of neon, we consider O VIII to be a better candidate for the observed emission line.
The relatively small redshift of the O VIII line ($`z`$ 0.01) can be interpreted in several ways, and we briefly discuss two possible scenarios. A value of $`z`$ = 0.008โ0.014 corresponds to the gravitational redshift of a photon emitted at a distance of 35โ60 Schwarzschild radii from a compact object: this interpretation would exclude the possibility that 4U 1700$`+`$24 is a white dwarf and would correspond to a distance of 1.5โ2.6$`\times `$10<sup>7</sup> cm from a 1.4 $`M_{}`$ neutron star. In a photoionized plasma, assuming that most of the oxygen is H-like, the measured O VIII line luminosity of $``$10<sup>40</sup> photons s<sup>-1</sup> gives an emission measure of $`EM6\times 10^{53}`$ cm<sup>-3</sup> (taking an abundance of $``$1.6$`\times 10^3`$, as derived from the oxygen edge fit and assuming that the overabundance is intrinsic to the Xโray source). Assuming a spherically symmetric geometry, if the O VIII line is emitted at $`2\times `$10<sup>7</sup> cm from the neutron star, we can estimate a density of $`n4\times 10^{15}`$ cm<sup>-3</sup>. For the measured continuum luminosity of 10<sup>34</sup> erg s<sup>-1</sup>, the ionization parameter is $`\xi =`$L$`{}_{cont}{}^{}/(`$n R$`{}_{}{}^{2})`$5000 erg cm s<sup>-1</sup>. For such a high ionization parameter, most of the oxygen should be fully ionized and therefore we would not expect to see a prominent O VIII emission line. O VIII dominates the ionization states of oxygen for $`\xi 100`$ erg cm s<sup>-1</sup>, which means that either we have used an oversimplified model (assuming, for example, that the line emission region is symmetric and homogeneous) or that the O VIII line is emitted at larger distance from the central X-ray source and therefore that the redshift is not solely gravitational.
In an alternative scenario, the same redshift can be produced by Doppler effects if the emitting plasma is moving away from us at a speed of 2000โ4000 km/s. This velocity is about two orders of magnitude larger than the wind velocity of an M-type giant (Reimers 1977), as well as the proper motion and possible orbital velocity of the binary system (Galloway, Sokoloski & Kenyon 2002). Although the high luminosity, hard X-ray spectrum, and rapid X-ray variability of 4U 1700$`+`$24 make it rather unlikely, the present data do not rule out the possibility that the accreting object is a white dwarf. The detection of Doppler-shifted lines of highly ionized elements has been reported for some supersoft sources and interpreted as the signature of collimated outflows (โjetsโ) coming from the accreting white dwarf (e.g., Cowley et al. 1998). Furthermore, the jet velocity is typically similar to the value we derive for the redshifted O VIII line. In our case, however, no corresponding blueshifted line is detected, which implies either a unipolar jet or some special geometry to obscure the approaching jet.
For a neutron star with a $`10^{12}`$ gauss magnetic field and a luminosity of $`10^{34}`$ erg s<sup>-1</sup>, the magnetospheric radius is $`3\times `$ 10<sup>9</sup> cm (Hayakawa 1985). Although X-ray pulsations have not been found from 4U 1700$`+`$24 , the fact that it is probably viewed close to face-on (Galloway et al. 2002) means that pulsations would be undetectable even if it is an accreting neutron star with high magnetic field as the line of sight is almost aligned with the neutron star spin axis. At a distance of $`3\times `$ 10<sup>9</sup> cm from a 1.4 $`M_{}`$ neutron star, a free particle would move radially towards the neutron star with a velocity of $`3400`$ km/s. Therefore, the redshift we measure could originate from close to the magnetospheric radius, where the plasma density is increased by the funneling effect produced by the magnetic field lines. Repeating the estimates of the ionization parameter described above, at a distance of 3$`\times `$ 10<sup>9</sup> cm the corresponding density is $`n2\times 10^{12}`$ cm<sup>-3</sup> and $`\xi 500`$. These two values are fairly consistent with the presence of an O VIII line.
Redshifted emission lines have also been reported in Cygnus X-3. Paerels et al. (2000) reported a redshift corresponding to 750โ800 km/s for all the emission lines detected in a Chandra HETGS observation of Cygnus X-3. In two more recent Chandra observations, a slightly smaller redshift of 270โ550 km/s, in addition to a Doppler modulation of about 150 km/s, was found (Stark & Saia 2003). The second effect was interpreted as being due to the binary orbital motion, but no convincing explanation was given for the global redshift. Although the X-ray luminosities and observed redshifts of 4U 1700$`+`$24 and Cygnus X-3 are rather different, a common mechanism could produce the line redshifts. For both 4U 1700$`+`$24 and Cygnus X-3, a low orbital inclination has been proposed (Galloway et al. 2002; Mioduszewski et al. 2001). Thus, Doppler shifts due to the bulk velocity of matter receding from us as it falls onto the compact object could be the origin of the redshifted lines.
###### Acknowledgements.
Based on observations obtained with XMM-Newton, an ESA science mission with instruments and contributions directly funded by ESA Member States and NASA. This work was partially supported by the Netherlands Organization for Scientific Research (NWO). This work was supported in part by the NASA Long Term Space Astrophysics program under grant NAG 5-9184. JMM and JLS gratefully acknowledge support from the NSF.
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# Reference potential approach to the quantum-mechanical inverse problem: I. Calculation of phase shift and Jost function
## 1. Introduction
Strict mathematical criteria for the unique solution of the inverse problem for the one-dimensional Schrรถdinger equation have been formulated more than 50 years ago. The history of this prestigious research area goes back to the dawn of quantum mechanics , and the decisive breakthrough has been achieved thanks to the important contributions by Borg , Levinson , Bargmann , Gelโfand and Levitan , Jost and Kohn , Marchenko , Krein , and others (see, e.g., for a thorough overview). As a result of fruitful brainstorming which culminated in early fifties it has been established that the quantum mechanical inverse problem can be solved if one manages to fix the so-called spectral function . In the case of confining one-dimensional quantum systems it means that the interaction potential can be uniquely determined if and only if the following complete set of information is available:
* Full energy spectrum of the bound states $`E_n`$ $`(n=0,1,\mathrm{},N).`$
* Full energy dependence (from 0 to $`\mathrm{}`$) of the phase shift $`\delta (E)`$ for the scattering states $`E>0`$.
* $`n`$ additional real parameters $`C_n`$ $`(n=0,1,\mathrm{},N)`$ related to relevant bound states that uniquely fix their normalization.
In principle, all these data can be obtained experimentally, but unfortunately, this is almost unachievable in practice. The real situation is even more hopeless, because in addition to the deficit of input information one inevitably faces very serious computational-technical problems. Thus, in spite of the whole mathematical beauty of the theory, one comes to a regrettable conclusion that rigorous solution of the quantum-mechanical inverse problem is a tremendously difficult task. For this reason most methods of deducing potentials from the available experimental data are based on some simplifying preconditions, therefore being inaccurate from the rigorous quantum mechanical point of view. Nevertheless, such methods may prove quite useful. For example, semiclassical approaches introduced many years ago by Rydberg , Klein , Rees and Dunham are still very popular in spectroscopy of diatomic molecules, and these concepts are constantly improved and developed .
Is it possible to apply rigorous methods of the inverse quantum theory for practical purposes? In view of the principle difficulties mentioned above, one has to be cautious in answering this question. In this paper we restrict ourselves to simple one-dimensional quantum systems, and the analysis proceeds from an idea that for any system of this kind one can build up a reasonable reference potential based on the available experimental data. As the reference potential is known, it is always possible to calculate all its discrete energy eigenvalues $`E_n`$, their norming constants $`C_n`$, and the phase shifts $`\delta (E)`$ for the scattering states. Therefore, in this artificial way one gets the full set of input information needed to uniquely solve the quantum mechanical inverse problem. Of course, the described approach is tautological: there is no need to regain a potential which is already known by definition. Nevertheless, such an approach is not meaningless, as it gives good zeroth approximations to the important spectral characteristics, such as Jost function and spectral density (the terms to be specified below). One can interpret the reference potential as only an initial guess to the real potential. Although the calculated quantities $`E_n`$, $`C_n`$ and $`\delta (E)`$ do not exactly match the actual values for the real system, they are still expected to be quite close to them. For a given reference potential one can calculate its Bargmann potential whose Jost function differs from the initial one only by a rational factor . By a suitable choice of this factor one can take a more adequate account of the experimental data. For example, one can replace the calculated discrete energy eigenvalues $`E_n`$ related to the reference potential with their actually observed values. Consequently, at least in some sense the new potential would be more realistic than the initial reference potential. There is also another motivation for the described โinverseโ approach to the inverse problem. Namely, through direct practical experience one can essentially increase his knowledge of how to overcome serious computational-technical difficulties when applying rigorous methods of the inverse quantum theory.
The one-dimensional inverse theory can be applied to diatomic molecules, since a two-particle problem can be always reduced to a one-particle problem in a spherically symmetric field. On this basis, we are going to examine the inverse problem for diatomic xenon molecule in its ground electronic state. Several reports of the research are planned. Methods of solution of the integral equations that enable to uniquely ascertain the potential are discussed and illustrated in the next paper of this series , while in this paper the emphasis is put on explaining the details of the basic concepts related to the reference potential approach. In Section 2 we briefly describe how the reference potential for the model system has been constructed, and how its discrete energy eigenvalues have been calculated. In Section 3 we describe the details of calculating the phase shift for the scattering states, and demonstrate full agreement with Levinson theorem . Section 4 aims to explain the important role of the Jost function in the quantum inverse theory. In particular, a detailed analysis of the asymptotic behavior of the Jost function is given. Finally, a brief conclusion is given in Section 5.
## 2. Exactly solvable reference potential for Xe<sub>2</sub>
In Fig. 1 one can see the reference potential constructed for the Xe<sub>2</sub> molecule. The same curve is depicted in both graphs, but very different energy scales are used. Throughout this paper only the rotationless case is analyzed, i.e., the rotational quantum number $`J=0.`$ According to the starting idea of the approach the only criterion for the choice of the reference potential is its agreement with the available experimental data. Therefore, we will not pay too much attention to various mathematical nuances and simply assume that the reference potential should be smooth and integrable in the whole physical domain. In addition, we try to construct a reference potential whose analytic form is as simple as possible. A good choice for this purpose, as explained in detail elsewhere , is a multi-component potential composed of smoothly joined Morse-type pieces
(1)
$$V(r)=V_k+D_k\left[\mathrm{exp}(\alpha _k(rr_k))1\right]^2,\text{ }r(0,\mathrm{}),$$
where $`V_k,`$ $`D_k,`$ $`\alpha _k`$ and $`r_k`$ are some constants (not definitely positive), and the subscript $`k`$ corresponds to different components smoothly joined at some suitably chosen boundary points $`X_{k+1}`$. The reference potential shown in Fig. 1 consists of only three components ($`k=0,1,2`$), the most internal of them ($`k=0`$) being a so-called pseudo-Morse potential. It means that the tiny potential well corresponding to this component (if taken separately) is just of the limit depth to entirely lose the discrete energy spectrum. Consequently, $`D_0=\mathrm{}^2\alpha _0^2/(8m)`$ ($`m`$ being the reduced mass of the pair of atoms), so that only three independent parameters remain for this component. The central component ($`k=1`$) is an ordinary Morse potential, while the most external one ($`k=2`$) is a โreversedโ Morse potential with the parameter $`D_2`$ being negative. By introducing a โreversedโ component one artificially creates a small potential hump in the long-distance range. This might seem unphysical and unjustified, but the point is that the height of this artificial hump approaches zero as the parameter $`r_2`$ approaches infinity. Therefore, taking a sufficiently large $`r_2`$, the hump becomes almost insignificant, while the analytic treatment remains simple and flexible. All parameters of the reference potential can be easily determined, if one requires continuity of the potential and its first derivative at the boundary points $`X_1`$ and $`X_2`$ (also shown in Fig. 1). These parameters as well as the calculated discrete energy eigenvalues $`E_n`$ $`(n=0,1,\mathrm{},23)`$ are given in Table 1.
The essence of the described construction is that the energy eigenvalue problem for the reference potential can be easily solved by solely analytic means to any desired accuracy . Moreover, as we demonstrate in the next Section, the major part of the energy dependence of the phase shift can also be ascertained analytically, which is a great advantage compared with applying numerical methods.
Most spectroscopic applications are related to the distance region shown in the lower graph of Fig. 1. For the inverse quantum theory, however, it is important to accurately reproduce the potential near the zero point $`r=0,`$ which is almost meaningless for spectroscopic applications. In this context one cannot ignore the fact that according to Eq. (1) the reference potential is finite at $`r=0`$ (see the upper graph of Fig. 1). Actually it means that the potential โjumpsโ to infinity at zero point, and this might also seem unphysical and unjustified. However, one has to bear in mind that the behavior of the real potential near $`r=0`$ is unknown and remains unknown. It does not matter so much how we describe the potential in this region, in so far as spectroscopic applications are of our main interest. As we see in the next section, description in terms of a pseudo-Morse potential is mathematically simple and elegant, and this, too, can be taken as a motivation for the approach to be used.
## 3. Phase shift and Levinson theorem
Next we are going to calculate the phase shift $`\delta (E)`$ for the full range of scattering states $`E(0,\mathrm{}).`$ To this end, one can use a long-known method first introduced by Morse and Allis (let us remind that only the rotationless case $`J=0`$ is examined here). It is based on solution of the following equation:
(2)
$$\delta ^{}(r,k)=\frac{\sqrt{2m}V(r)}{\mathrm{}k}\mathrm{sin}^2\left[kr+\delta (r,k)\right],\text{ }\delta (0,k)=0,$$
where $`k=`$ $`{\displaystyle \frac{\sqrt{2mE}}{\mathrm{}}}.`$ The phase shift is then determined as $`\delta (k)=\underset{r\mathrm{}}{lim}\delta (r,k).`$ The described method is universal but rather time-consuming, because for any energy from 0 to $`\mathrm{}`$ one has to perform an integration from 0 to $`\mathrm{}.`$ Fortunately, at this point we can take advantage of the special analytic form of the reference potential. As mentioned, the region $`r(0,X_1)`$ is approximated by a pseudo-Morse potential. It means that one immediately gets two linearly independent analytic solutions of the corresponding Schrรถdinger equation
(3) $`\mathrm{\Psi }_0^{(1)}(r)`$ $`=\mathrm{exp}(y_0/2)y_0^{i\mu _0}\mathrm{\Psi }(i\mu _0,2i\mu _0+1;y_0)`$
$`\mathrm{\Psi }_0^{(2)}(r)`$ $`=\mathrm{exp}(y_0/2)(y_0)^{i\mu _0}\mathrm{\Psi }(i\mu _0,2i\mu _0+1;y_0),\text{ }r(0,X_1),`$
where $`\mu _0=1/2\sqrt{(EV_0)/D_01}>0`$, $`y_0\mathrm{exp}\left[\alpha _0(rr_0)\right]`$, and $`\mathrm{\Psi }(i\mu _0,2i\mu _0+1;y_0)`$ is a particular solution of the confluent hypergeometric equation introduced by Tricomi (see for details). If $`\mu _0^2<<y_0,`$ this function can be evaluated from the asymptotic series
(4)
$$\mathrm{\Psi }(i\mu _0,2i\mu _0+1;y_0)=y_0^{i\mu _0}\underset{n=0}{\overset{N}{}}\frac{(i\mu _0)_n(i\mu _0)_n}{n!(y_0)^n},$$
where $`(a)_n\mathrm{\Gamma }(a+n)/\mathrm{\Gamma }(a)=a(a+1)(a+2)\mathrm{}(a+n1)`$ is the Pochhammer symbol, and $`N`$ must not be too large. Thus
(5)
$$\mathrm{\Psi }_0^{(1)}(r)=\mathrm{exp}(y_0/2)\left[1\frac{\mu _0^2}{1!y_0}+\frac{\mu _0^2(\mu _0^2+1^2)}{2!y_0^2}\frac{\mu _0^2(\mu _0^2+1^2)(\mu _0^2+2^2)}{3!y_0^3}+\mathrm{}\right],$$
and analogously
(6)
$$\mathrm{\Psi }_0^{(2)}(r)=\mathrm{exp}(y_0/2)\left[1+\frac{\mu _0^2}{1!y_0}+\frac{\mu _0^2(\mu _0^2+1^2)}{2!y_0^2}+\frac{\mu _0^2(\mu _0^2+1^2)(\mu _0^2+2^2)}{3!y_0^3}+\mathrm{}\right].$$
The phase shift is related to
regular solutions of the Schrรถdinger equation, which means that the physically correct linear combination of$`\mathrm{\Psi }_0^{(1)}(r)`$ and $`\mathrm{\Psi }_0^{(2)}(r)`$ should vanish as $`r0`$, i.e.,
$$\mathrm{\Psi }_0(r)=N_1\mathrm{\Psi }_0^{(1)}(r)+N_2\mathrm{\Psi }_0^{(2)}(r),$$
where
(7)
$$\frac{N_2}{N_1}=\mathrm{exp}\left[y_0(0)\right]\frac{1{\displaystyle \frac{\mu _0^2}{1!y_0(0)}}+{\displaystyle \frac{\mu _0^2(\mu _0^2+1^2)}{2!y_0^2(0)}}\mathrm{}}{1+{\displaystyle \frac{\mu _0^2}{1!y_0(0)}}+{\displaystyle \frac{\mu _0^2(\mu _0^2+1^2)}{2!y_0^2(0)}}+\mathrm{}}.$$
From Eq. (7) one can see that the particular solution $`\mathrm{\Psi }_0^{(2)}(r)`$ practically does not contribute to the regular solution at distances sufficiently far from zero point, but still deep inside the classically forbidden region. Indeed, in this case
(8)
$$\frac{N_2\mathrm{\Psi }_0^{(2)}(r)}{N_1\mathrm{\Psi }_0^{(1)}(r)}=\mathrm{exp}[y_0(r)y_0(0)]\times \frac{1{\displaystyle \frac{\mu _0^2}{1!y_0(0)}}+{\displaystyle \frac{\mu _0^2(\mu _0^2+1^2)}{2!y_0^2(0)}}\mathrm{}}{1+{\displaystyle \frac{\mu _0^2}{1!y_0(0)}}+{\displaystyle \frac{\mu _0^2(\mu _0^2+1^2)}{2!y_0^2(0)}}+\mathrm{}}\times $$
$$\times \frac{1+{\displaystyle \frac{\mu _0^2}{1!y_0(r)}}+{\displaystyle \frac{\mu _0^2(\mu _0^2+1^2)}{2!y_0^2(r)}}+\mathrm{}}{1{\displaystyle \frac{\mu _0^2}{1!y_0(r)}}+{\displaystyle \frac{\mu _0^2(\mu _0^2+1^2)}{2!y_0^2(r)}}\mathrm{}}=\mathrm{exp}[y_0(r)y_0(0)]\times $$
$$\times \frac{1+{\displaystyle \frac{\mu _0^2}{1!y_0(0)}}\left[\mathrm{exp}(\alpha _0r)1\right]+\mathrm{}}{1{\displaystyle \frac{\mu _0^2}{1!y_0(0)}}\left[\mathrm{exp}(\alpha _0r)1\right]+\mathrm{}}\mathrm{exp}\{\mathrm{exp}(\alpha _0r_0)[1\mathrm{exp}(\alpha _0r)]\}$$
is an extremely small quantity.
There is still a lot more profit to gain from the pseudo-Morse approximation to calculate the phase shift. As we just proved, the physically correct solution (apart from normalization) in a wide energy range (practically up to $`EV(0)`$) reduces to the particular solution $`\mathrm{\Psi }_0^{(1)}(r).`$ Using some well-known relations from the theory of confluent hypergeometric functions this solution can be rewritten
(9)
$$\mathrm{\Psi }_0^{(1)}(r)=\mathrm{exp}(y_0/2)y_0^{i\mu _0}\mathrm{\Psi }(i\mu _0,2i\mu _0+1;y_0)=$$
$$\mathrm{exp}(y_0/2)\left[\frac{\mathrm{\Gamma }(2i\mu _0)}{\mathrm{\Gamma }(i\mu _0)}y_0^{i\mu _0}\mathrm{\Phi }(i\mu _0,2i\mu _0+1;y_0)+\frac{\mathrm{\Gamma }(2i\mu _0)}{\mathrm{\Gamma }(i\mu _0)}y_0^{i\mu _0}\mathrm{\Phi }(i\mu _0,2i\mu _0+1;y_0)\right],$$
where the symbols $`\mathrm{\Phi }(a,c;x)=1+{\displaystyle \frac{a}{c}}{\displaystyle \frac{x}{1!}}+{\displaystyle \frac{a(a+1)}{c(c+1)}}{\displaystyle \frac{x^2}{2!}}+\mathrm{}`$ denote the well-known confluent hypergeometric functions.
As is seen, Eq. (9) represents a sum of two complex conjugates. Therefore,
$$\mathrm{\Psi }_0^{(1)}(r)A_0(y_0)\mathrm{cos}\left[B_0(y_0)\phi _0\alpha _0\mu _0r\right],$$
where the functions $`A_0`$ and $`B_0`$ can be determined from a series :
$`A_0(y_0)e^{iB_0(y_0)}`$ $`=1{\displaystyle \frac{y_0/4}{i\mu _0+1/2}}+{\displaystyle \frac{\left(y_0/4\right)^2}{\left(i\mu _0+1/2\right)1!}}\left(1{\displaystyle \frac{y_0/4}{i\mu _0+3/2}}\right)+`$
(10) $`+`$ $`{\displaystyle \frac{\left(y_0/4\right)^4}{\left(i\mu _0+1/2\right)\left(i\mu _0+3/2\right)2!}}\left(1{\displaystyle \frac{y_0/4}{i\mu _0+5/2}}\right)+\mathrm{}`$
A good point is that the phase parameter (which, of course, is not the actual phase shift) $`\phi _0`$ $`\alpha _0\mu _0r_0\mathrm{arg}\left[{\displaystyle \frac{\mathrm{\Gamma }(2i\mu _0)}{\mathrm{\Gamma }(i\mu _0)}}\right]`$ can be calculated exactly. Indeed, using a Legendre formula for doubling a gamma functionโs argument
$$\mathrm{\Gamma }(2z)=\frac{2^{2z1}}{\sqrt{\pi }}\mathrm{\Gamma }(z)\mathrm{\Gamma }(z+1/2),$$
and another useful formula
(11)
$$\mathrm{arg}\mathrm{\Gamma }(i\mu _0+1/2)=\mu _0\left(\frac{1}{2}\mathrm{ln}(1+4\mu _0^2)1\mathrm{ln}2\right)\frac{1}{2}\underset{0}{\overset{\mathrm{}}{}}\left(\mathrm{coth}t\frac{1}{t}\right)e^t\mathrm{sin}(2\mu _0t)\frac{dt}{t},$$
one gets the following result:
(12)
$$\phi _0=\mu _0\left(\alpha _0r_0+1\mathrm{ln}2\frac{1}{2}\mathrm{ln}(1+4\mu _0^2)\right)+\frac{1}{2}\underset{0}{\overset{\mathrm{}}{}}\left(\mathrm{coth}t\frac{1}{t}\right)e^t\mathrm{sin}(2\mu _0t)\frac{dt}{t},$$
where the integral can be conveniently evaluated
(13)
$$I\underset{0}{\overset{\mathrm{}}{}}\left(\mathrm{coth}t\frac{1}{t}\right)e^t\mathrm{sin}(2\mu _0t)\frac{dt}{t}=\underset{0}{\overset{T}{}}e^t\mathrm{sin}(\pi \frac{t}{T})f(t)๐t,\text{ }T=\frac{\pi }{2\mu _0},$$
$$f(t)=\frac{\mathrm{coth}t\frac{1}{t}}{t}e^T\frac{\mathrm{coth}(t+T)\frac{1}{t+T}}{t+T}+e^{2T}\frac{\mathrm{coth}(t+2T)\frac{1}{t+2T}}{t+2T}\mathrm{}$$
Another equivalent formula for this quantity reads
(14)
$$I=\underset{n=1}{\overset{\mathrm{}}{}}I_n,\text{ }I_n=\frac{(1)^{n1}2^{2n}B_n}{(2n)(2n1)(1+4\mu _0^2)^{2n1}}\underset{k=0}{\overset{n1}{}}(1)^k\left(\genfrac{}{}{0pt}{}{2n1}{2k+1}\right)(2\mu _0)^{2k+1}.$$
Here $`B_n`$ denotes the $`n`$-th order Bernoulli number.
Correct linear combinations of solutions of the Schrรถdinger equations for analytically different pieces of the reference potential are uniquely fixed by the continuity requirements of the wave function and its first derivative at the boundary points $`X_1`$ and $`X_2`$. Since the parameter $`\phi _0`$ can be calculated exactly, the phase shift can also be ascertained with the help of solely analytic means in the energy range $`0<EV(0).`$ At high energies ($`EV(0)`$), however, the analytic approach gradually fails, since the particular solution $`\mathrm{\Psi }_0^{(2)}(r)`$ in Eq. (3) cannot be ignored any more. To determine the phase shift in this region one has to solve Eq. (2). Fortunately, there is no need to perform integration over the whole physical domain $`r(0,\mathrm{})`$, because the solution in the long-distance range $`rX_2`$ already has the โrightโ analytic form
(15)
$$\mathrm{\Psi }_2(r)=2A_2(y_2)\mathrm{cos}\left[B_2(y_2)+\phi _2kr\right]=2A_2(y_2)\mathrm{sin}\left[\delta (r,k)+kr\right],$$
from which the phase shift can be easily obtained. Here $`y_2(r)2a_2\mathrm{exp}\left[\alpha _2(rr_2)\right],`$ $`a_2\sqrt{2mD_2}/\left(\mathrm{}\alpha _2\right),`$ $`\delta (r,k)`$ is the solution of Eq. (2), and the complex function $`A_2(y_2)\mathrm{exp}\left[iB_2(y_2)\right]\mathrm{exp}\left(iy_2/2\right)\mathrm{\Phi }[i(k/\alpha _2a_2)+1/2,2ik/\alpha _2+1;iy_2].`$ Since $`A_2(y_2)0`$ and $`B_2(y_2)1`$ as $`r0,`$ the phase shift
(16)
$$\delta (k)=\delta (X_2,k)+B_2\left[y_2(X_2)\right],$$
i.e., one only needs to integrate Eq. (2) until the boundary point $`X_2`$.
Now, let us recall an important relation known as Levinson theorem
(17)
$$\delta (0)\delta (\mathrm{})=n\pi ,$$
which correlates the energy dependence of the phase shift with the number of bound states. As can be seen in Fig. 2, a really good agreement with the Levinson theorem can be obtained, but only if the phase shift is calculated up to very high energies (note that the energy scale in Fig. 2 is logarithmic, and it involves 20 orders of magnitude!). Moreover, the phase shift has to be calculated with sufficiently high precision throughout the whole energy range, otherwise there is no chance to accurately ascertain other important spectral characteristics, such as Jost function. Unfortunately, the higher the energy goes, the more complicated and time-consuming the numerical integration of Eq. (2) becomes. How could we bridge over this troublesome technical difficulty? In such situation, one may recall some general principles, and this is indeed helpful to complete calculations of the phase shift.
Let us express the phase shift as formal result of integration of Eq. (2):
(18)
$$\delta (k)=\underset{r\mathrm{}}{lim}\delta (r,k)=\frac{1}{2Ck}\underset{0}{\overset{\mathrm{}}{}}V(r)\left\{1\mathrm{cos}\left[2kr+2\delta (r,k)\right]\right\}๐r=\frac{\underset{0}{\overset{\mathrm{}}{}}V(r)๐r}{2Ck}+$$
$$+\frac{1}{2Ck}\underset{0}{\overset{\mathrm{}}{}}\mathrm{cos}\left(2kr\right)\left\{V(r)\mathrm{cos}\left[2\delta (r,k)\right]\right\}๐r\frac{1}{2Ck}\underset{0}{\overset{\mathrm{}}{}}\mathrm{sin}\left(2kr\right)\left\{V(r)\mathrm{sin}\left[2\delta (r,k)\right]\right\}๐r.$$
Here a special denotation has been introduced (and will be used henceforward) for the constant $`C{\displaystyle \frac{\mathrm{}^2}{2m}}`$ that often appears in formulas. Now, let us call for help from the famous Riemann-Lebegue theorem (see, e.g., ): if a function $`F(r)`$ is integrable in an interval $`r(a,b)`$, then
(19)
$$\underset{a}{\overset{b}{}}\mathrm{cos}\left(\lambda r\right)F(r)๐r0,\text{ }\underset{a}{\overset{b}{}}\mathrm{sin}\left(\lambda r\right)F(r)๐r0\text{ as }\lambda \mathrm{}.$$
From this one immediately concludes that the last two integrals in Eq. (18) will vanish as $`k\mathrm{}.`$ However, we can integrate by parts and apply the Riemann-Lebegue theorem to the resulting integrals. This procedure can be repeated as many times as needed, and as a result, one comes to a rather general asymptotic formula for the phase shift (it can easily proved that only odd powers of $`k`$ appear in this series)
(20)
$$\delta (k)=\frac{a_1}{k}+\frac{a_3}{k^3}+\frac{a_5}{k^5}+\mathrm{},\text{ }k\mathrm{},$$
where
(21)
$$a_1=\frac{\underset{0}{\overset{\mathrm{}}{}}V(r)๐r}{2Ck},\text{ }a_3=\frac{1}{8C^2}\left\{\underset{0}{\overset{\mathrm{}}{}}\left[V(r)\right]^2๐r+CV^{}(0)\right\},\text{ }$$
$$a_5=\frac{1}{32C^3}\left\{\underset{0}{\overset{\mathrm{}}{}}\left[V(r)\right]^3๐r+C\underset{0}{\overset{\mathrm{}}{}}\left[V^{}(r)\right]^2๐r+4CV(0)V^{}(0)C^2V^{\prime \prime \prime }(0)\right\}.$$
Thus one can easily calculate the coefficients $`a_1`$, $`a_3`$, $`a_5`$, etc., and this is just what is needed to ascertain the whole energy dependence of the phase shift. Fig. 3 demonstrates how well Eq. (20) fits with the results of direct numerical integration of the phase equation.
One cannot so easily find any direct illustrations to the Levinson theorem from the literature (at least the author of this paper has not found them), and Fig. 2, which is an illustration of this kind, could therefore be of more general interest than merely an attachment to a particular model potential. The curve shown in this figure has been calculated with at least 8 significant digits, and it is in full agreement with all relevant general physical considerations. This concerns not only the Levinson theorem and the asymptotic behavior at large energies, but also the low energy part of the energy dependence (see the left-side inset of Fig. 2), exactly corresponding to the well-known formula
(22)
$$\delta (k)=n\pi \mathrm{arctan}(kR_0),\text{ }k0$$
($`R_0`$ being scattering length), which can be found in most handbooks on quantum mechanics.
## 4. Jost function and inverse problem
Having calculated the phase shift, we have come to a situation from which an ideal inverse problem study would start. In other words, we are now provided with the full set of information needed to uniquely solve the inverse problem. In our case this would mean that we simply regain the reference potential from which we started. This, of course, is not our main goal. As explained in Section 1, we are interested in getting realistic zeroth approximations to the important spectral functions which then could be used to improve the initial model potential. The latter step, however, is planned as a subject for a forthcoming publication. In this paper we only try to make sure that the described scheme is reliable, and to this end the next step is to accurately calculate the Jost function, the main spectral characteristic, which contains the most part of information needed to solve the inverse problem. A thorough overview of all useful properties of the Jost function is given in Ref. . For the treatment here, the most important point is that the Jost function creates a simple link between physical and regular solutions of the scattering states. The physical solution ($`J=0`$) reads (cf. with Eq. (15))
(23)
$$\mathrm{\Psi }(r,k)\mathrm{exp}i\delta (k)\mathrm{sin}\left[kr+\delta (k)\right],\text{ }r\mathrm{},$$
while the regular solution $`\phi (r,k)`$ is defined by a condition
(24)
$$\underset{r0}{lim}\frac{\phi (r,k)}{r}=1.$$
It can be shown that these two solutions of the Schrรถdinger equation are proportional
(25)
$$\mathrm{\Psi }(r,k)=\frac{k}{F(k)}\phi (r,k),$$
where
(26)
$$F(k)=\left|F(k)\right|\mathrm{exp}\left[i\delta (k)\right]$$
is the Jost function we are talking about. For further treatment we have to calculate the modulus of this function
(27)
$$\left|F(E)\right|=\underset{n=0}{\overset{N}{}}(1E_n/E)\mathrm{exp}\left[\frac{1}{\pi }P\underset{0}{\overset{\mathrm{}}{}}\frac{\delta (E^{})dE}{E^{}E}\right],\text{ }E(0,\mathrm{}).$$
Here $`E_n`$ are the discrete energy eigenvalues and the symbol $`P`$ points at the principal value of the integral. As the phase shift and the bound states are known, the calculations are relatively trivial, but they have to be carried out very accurately. Next we can fix the spectral density
(28)
$$\frac{d\rho (E)}{dE}=\{\genfrac{}{}{0.0pt}{}{\pi ^1\sqrt{E}|F(E)|^2,E0,}{\underset{n}{}C_n\delta (EE_n),\text{ }E<0.}$$
Here $`C_n`$ are the norming constants for the relevant bound states. Note that these quantities are related to regular solutions, wherefore their ascertainment is not so easy task as one might think.
Now we have come close to the real solution schemes of the inverse problem. For example, the Gelfand-Levitan method is based on the integral equation
(29)
$$K(r,r^{})+G(r,r^{})+\underset{0}{\overset{r}{}}K(r,s)G(s,r^{})๐s=0,$$
whose kernel reads
(30)
$$G(r,r^{})=\underset{\mathrm{}}{\overset{\mathrm{}}{}}\frac{\mathrm{sin}\left(kr\right)\mathrm{sin}\left(kr^{}\right)}{k^2}๐\sigma ,$$
and the quantity $`d\sigma d\rho (E){\displaystyle \frac{d\rho _0(E)}{dE}}dE.`$ Here $`{\displaystyle \frac{d\rho _0(E)}{dE}}`$ is free particleโs spectral density (related to the potential $`V(r)0`$), and therefore,
(31)
$$d\sigma =\{\genfrac{}{}{0.0pt}{}{d\rho (E)d({\displaystyle \frac{2E^{3/2}}{3\pi }}),\text{ }E0}{d\rho (E),\text{ }E<0.}$$
Eq. (30) can be rewritten
(32)
$$G(r,r^{})=\frac{2}{\pi }\underset{0}{\overset{\mathrm{}}{}}\mathrm{sin}\left(kr\right)\mathrm{sin}\left(kr^{}\right)g(k)๐k+\underset{n}{}\frac{C_n}{4\gamma _n^2}\mathrm{sinh}\left(\gamma _nr\right)\mathrm{sinh}\left(\gamma _nr^{}\right),$$
where $`\gamma _n^2={\displaystyle \frac{2mE_n}{\mathrm{}^2}}`$ and the function
(33)
$$g(k)\frac{1}{\left|F(k)\right|^2}1.$$
If one is able to solve Eq. (29), the potential can be determined from the relation
(34)
$$V(r)=2C\frac{d}{dr}K(r,r).\text{ }$$
Eqs. (29) to (34) explicitly demonstrate great importance of the Jost function in the inverse quantum theory. The energy dependence of the Jost functionโs modulus is shown in Fig. 4. One can see that at small energies this quantity achieves extremely large values. Near the โcriticalโ energy $`E=V(0)`$ the curve rapidly turns from nearly vertical to nearly horizontal, and at still higher energies it slowly approaches the limit value $`\left|F(k)\right|1`$, โbreaking freeโ from the potential field (note that for a free particle, $`F(k)=1`$ independent of energy).
The kernel of the Gelfand-Levitan equation is essentially determined by the function $`g(k)`$ given by Eq. (33), which means that this function has to be ascertained very accurately, in order to solve Eq. (29) and calculate the potential according to Eq. (34). As can be seen in Fig. 5, $`g(k)`$ noticeably differs from unity only at $`kk_0\sqrt{V(0)/C}`$, and $`g(k)0`$ as $`k\mathrm{}.`$ Naturally, as shown in the inset, there is no break of derivative in the โcriticalโ region. Since the asymptotic behavior of the function $`g(k)`$ essentially determines the shape of the potential near the zero point $`r=0`$ (see the end of this section), it makes sense to analyze this behavior in more details. We have already ascertained the asymptotic expression for the phase shift (see Eq. (20)), and this can be used to immediately get the asymptotic formulas for both the Jost functionโs modulus and the function $`g(k)`$. For example, taking
(35)
$$\mathrm{ln}\left|F(k)\right|=\frac{a_2}{k^2}+\frac{a_4}{k^4}+\frac{a_6}{k^6}+\mathrm{},\text{ }k\mathrm{}$$
(this time only even powers of $`k`$ appear in the series, as can be easily proved), the coefficients read
(36) $`a_2`$ $`={\displaystyle \frac{2}{\pi }}(a_1k_a{\displaystyle \frac{a_3}{k_a}}{\displaystyle \frac{a_5}{3k_a^3}}\mathrm{})+{\displaystyle \frac{1}{C}}\left[{\displaystyle \frac{1}{\pi }}{\displaystyle \underset{0}{\overset{E_a}{}}}\delta (E^{})๐E^{}{\displaystyle \underset{n}{}}E_n\right],`$
$`a_4`$ $`={\displaystyle \frac{2}{\pi }}({\displaystyle \frac{a_1k_a^3}{3}}+a_3k_a{\displaystyle \frac{a_5}{k_a}}\mathrm{})+{\displaystyle \frac{1}{C^2}}\left[{\displaystyle \frac{1}{\pi }}{\displaystyle \underset{0}{\overset{E_a}{}}}\delta (E^{})E^{}๐E^{}{\displaystyle \frac{1}{2}}{\displaystyle \underset{n}{}}\left(E_n\right)^2\right],`$
$`a_6`$ $`={\displaystyle \frac{2}{\pi }}({\displaystyle \frac{a_1k_a^5}{5}}+{\displaystyle \frac{a_3k_a^3}{3}}+a_5k_a\mathrm{})+{\displaystyle \frac{1}{C^3}}\left[{\displaystyle \frac{1}{\pi }}{\displaystyle \underset{0}{\overset{E_a}{}}}\delta (E^{})\left(E^{}\right)^2๐E^{}{\displaystyle \frac{1}{3}}{\displaystyle \underset{n}{}}\left(E_n\right)^3\right],`$
where $`E_a=Ck_a^2`$ is an arbitrary energy value in the range where the asymptotic approximation Eq. (20) can be used.
Eq. (36) may look nice but it is a bit inconvenient in practice. Fortunately, a straightforward approach exists, which enables to ascertain the coefficients in Eq. (35) more accurately and much more easily without any direct reference to the phase shift. The approach in question is based on the following integral representation for the Jost function :
(37)
$$F(k)=1+\frac{1}{C}\underset{0}{\overset{\mathrm{}}{}}e^{ikr}V(r)\phi (k,r)๐r,$$
where the regular solution $`\phi (k,r)`$ can be calculated with the help of a well-known iteration method. Namely, taking $`\phi ^{(0)}(k,r)={\displaystyle \frac{\mathrm{sin}kr}{k}}`$, and
(38)
$$\phi ^{(n)}(k,r)=\frac{1}{C}\underset{0}{\overset{r}{}}\frac{\mathrm{sin}k(rr^{})}{k}V(r^{})\phi ^{(n1)}(k,r^{})๐r^{},\text{ }n=1,2,\mathrm{},$$
the desired solution reads
(39)
$$\phi (k,r)=\underset{n=0}{\overset{\mathrm{}}{}}\phi ^{(n)}(k,r).$$
In Eq. (38) one can use integration by parts to calculate step-by-step the terms $`\phi ^{(1)}(k,r),`$ $`\phi ^{(2)}(k,r),`$ $`\phi ^{(3)}(k,r),\mathrm{}`$, and their contributions to the Jost function. Let us see, how to ascertain the correct asymptotic formula for the function $`g(k)={\displaystyle \frac{1}{\left|F(k)\right|^2}}1`$ as $`k\mathrm{},`$ which would include the terms until $`1/k^4.`$ Within this approximation
(40) $`\phi ^{(1)}(k,r)`$ $`={\displaystyle \frac{\mathrm{cos}kr}{2Ck^2}}W(r)+{\displaystyle \frac{\mathrm{sin}kr}{4Ck^3}}\left[V(0)+V(r)\right]+{\displaystyle \frac{\mathrm{cos}kr}{8Ck^4}}\left[V^{}(r)V^{}(0)\right],\text{ }W(r){\displaystyle \underset{0}{\overset{r}{}}}V(r^{})๐r^{},`$
$`\phi ^{(2)}(k,r)`$ $`={\displaystyle \frac{\mathrm{sin}kr\left[W(r)\right]^2}{8C^2k^2}}{\displaystyle \frac{\mathrm{cos}kr}{8C^2k^4}}\left[V(0)W(r)+V(r)W(r)+U(r)\right],\text{ }U(r){\displaystyle \underset{0}{\overset{r}{}}}\left[V(r^{})\right]^2๐r^{},`$
$`\phi ^{(3)}(k,r)`$ $`={\displaystyle \frac{\mathrm{cos}kr\left[W(r)\right]^3}{48C^3k^4}}`$
(all higher order terms can be ignored). Thereafter, using integration by parts in Eq. (37), one comes to the following formulas for the real and imaginary parts of the Jost function:
(41) $`\mathrm{Re}F(k)`$ $`=1+{\displaystyle \frac{V(0)}{4Ck^2}}{\displaystyle \frac{W^2}{8C^2k^2}}{\displaystyle \frac{V^{\prime \prime }(0)}{16Ck^4}}+{\displaystyle \frac{1}{32C^2k^4}}\left\{5\left[V(0)\right]^22V^{}(0)W\right\}`$
$`{\displaystyle \frac{1}{32C^3k^4}}\left[V(0)W^2+2UW\right]+{\displaystyle \frac{W^4}{384C^4k^4}},`$
(42)
$$\mathrm{Im}F(k)=\frac{W}{2Ck}+\frac{V^{}(0)}{8Ck^3}+\frac{1}{8C^2k^3}\left[V(0)W+U\right]\frac{W^3}{48C^3k^3},$$
where $`W\underset{0}{\overset{\mathrm{}}{}}V(r^{})๐r^{}`$ and $`U\underset{0}{\overset{\mathrm{}}{}}\left[V(r^{})\right]^2๐r^{}.`$ Quite surprisingly, when calculating $`\left|F(k)\right|^2=\left[\mathrm{Re}F(k)\right]^2+\left[\mathrm{Im}F(k)\right]^2`$, all troublesome terms will cancel out, resulting in nice expressions for both quantities of interest:
(43)
$$\left|F(k)\right|^2=1+\frac{V(0)}{2Ck^2}+\frac{3\left[V(0)\right]^2}{8C^2k^4}\frac{V^{\prime \prime }(0)}{8Ck^4},$$
(44)
$$g(E)=\frac{V(0)}{2E}+\frac{\left[V(0)\right]^2CV^{\prime \prime }(0)}{8E^2}.$$
Comparing Eqs. (35) and (44), one finds that the coefficients
(45)
$$a_2=\frac{V(0)}{4C},\text{ }a_4=\frac{2\left[V(0)\right]^2CV^{\prime \prime }(0)}{16C^2}$$
do not depend on the phase shift, but are directly related to the potential and its second derivative at the zero point, both these quantities being finite according to the starting idea of the approach. Thus we have described all properties of the Jost function in the frame of the proposed approach, and are now prepared to start calculation of the potential itself.
## 5. Conclusion
In this paper we proposed as if an โinversedโ approach to the quantum mechanical inverse problem. Starting from a known reference potential, one can calculate the important spectral characteristics of the system, the phase shift and the Jost function, which are almost unattainable in a real experiment. On one hand, the reference potential has to be realistic enough to be used as a zeroth approximation to the real potential. Only in this case there is a chance to construct a Bargmann potential that would be even more realistic, for example, whose discrete energy levels would exactly fit with the actually observed ones. On the other hand, we suggest to choose a reference potential that would be exactly solvable, in the sense that its energy eigenvalue problem can be solved to any desired accuracy with the help of solely analytic means. To this end, as demonstrated in Sections 2 and 3, a multi-component potential composed of smoothly joined Morse-type pieces is especially suitable. In particular, we would like to highlight the usefulness of the pseudo-Morse approximation for the small distances (and high energies) region of the potential. In this paper we only used a single pseudo-Morse component, which stretches until the zero point $`r=0`$. As demonstrated elsewhere , adding more pseudo-Morse components does not bring along any serious problems, so one can include just as many components of this type as he considers reasonable.
The special analytic form of the reference potential enables to ascertain the phase shift analytically up to the energies $`EV(0)`$ (note that an arbitrarily large value for $`V(0)`$ can be taken, if one introduces a sufficient number of pseudo-Morse components). As the asymptotic behavior of the phase shift can be determined from general physical considerations, it is possible to directly demonstrate full agreement with the famous Levinson theorem , as well as to ascertain the full energy dependence of the Jost function. Provided with this important input information, one can attack the real computational-technical problem, trying to solve the main integral equation in the frame of the Gelfand-Levitan , Marchenko or Krein method. This will be the subject of the next paper in this series .
## Acknowledgement
The research described in this paper has been supported by Grants No 5863 and 4508 from the Estonian Science Foundation.
## Figure captions
1. Three-component model potential for the system (Xe<sub>2</sub> in ground electronic state) investigated. Note that the same potential curve is depicted in both graphs (the lower one starts where the upper one ends), but essentially different energy scales are used for them. All components have the well-known analytic form of the Morse potential, but the ordinary Morse approximation is used only in the central range $`r[X_1,X_2]`$ (see the explanations in Section 2). The parameters of the components as well as the calculated discrete energy levels (24 in total) are given in Table 1.
2. Direct demonstration of the Levinson theorem ($`\delta (0)\delta (\mathrm{})=n\pi `$) for the model system studied. As needed, $`\delta (0)=24\pi `$, since the system has 24 bound states. At $`E=`$ 3.146294 meV, the phase shift passes a zero, and then remains negative, very slowly approaching the limit ($`\delta (\mathrm{})=0`$) as $`E\mathrm{}.`$ The left side-inset shows the nearly linear energy dependence as $`E0`$, in full agreement with Eq. (22). In the right-side inset one can see that the phase curve has an inflection point near $`E=V(0)`$.
3. A demonstration of the exellent agreement between the numerically calculated phase shifts and the general asymptotic formula, Eq. (20), as $`E\mathrm{}`$.
4. Demonstration of the results of calculating the Jost function (in fact, its modulus) for the reference potential. Again, as in the case of Fig. 1, different parts of the same curve are shown in the graphs and their insets, but rather different scales are used (note that in some cases the scale is logarithmic). The information gathered into the upper graph would be useful, if oneโs aim is to construct another potential whose Jost function only slightly differs from the initial one. The lower graph demonstrates how abruptly the curve turns from practically vertical to nearly horizontal in the vicinity of the โcriticalโ energy value $`E=V(0)`$, while this dramatic change is absolutely invisible in the scale used for the ordinate axis in the upper graph. The asymptotic formula Eq. (35) is used to calculate the Jost function for $`kk_a=75000`$ ร
(which corresponds to about 1.8\*10<sup>5</sup> eV) with coefficients given by Eq. (45).
5. Visualization of the characteristic function $`g(k)`$ (given gy Eq. (33)), which determines the kernel of the Gelfand-Levitan integral equation. As can be seen, $`g(k)=1`$ (with very high precision) if $`kk_0=\sqrt{V(0)/C}`$ (see Fig. 4 to undestand the reason for this), then starts to decrease, gradually โbreaking freeโ from the potential field, and $`g(k)0`$ as $`k\mathrm{}`$. In spite of the seeming simplicity of the energy dependence, the whole curve has to be (and has been) calculated very accurately up to very high energies, in order to accurately ascertain the potential. The high-energy part of the curve ($`kk_a`$) has been calculated according to Eq. (44).
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# Virtual QCD corrections to Higgs boson plus four parton processes
## I Introduction
In this paper we study the production of a standard model Higgs boson in association with two jets. This is one of the most promising discovery channels at the LHC especially for a Higgs boson with a mass in the range $`110`$ GeV$`<M_H<180`$ GeV. At Born level there are two classes of processes which contribute, as illustrated in Fig. 1.
In Fig. 1(a), the Higgs is produced via vector boson fusion, while in Fig. 1(b) the coupling of the Higgs boson to gluons is mediated by a top quark loop. In the limit in which the mass of the top quark tends to infinity the coupling can be treated using an effective theory as described below. We shall refer to this process as the gluon fusion process. Notice that the external gluons lines in Fig. 1(b) could as well be replaced by quarks.
The final aim of this study is the calculation of the Higgs + $`2`$-jet rate, at next-to-leading (NLO) order, where the Higgs is produced using the effective coupling to gluons,
$$_{\mathrm{eff}}=\frac{1}{4}A(1+\mathrm{\Delta })HG_{\mu \nu }^aG^{a\mu \nu }.$$
(1)
In Eq. (1), $`G_{\mu \nu }^a`$ is the field strength of the gluon field and $`H`$ is the Higgs-boson field. The effective coupling $`A`$ is given by
$$A=\frac{g^2}{12\pi ^2v},$$
(2)
where $`g`$ is the bare strong coupling and $`v`$ is the vacuum expectation value parameter, $`v^2=(G_F\sqrt{2})^1=(246\mathrm{GeV})^2`$. The finite $`O(g^2`$) correction to the effective operator has been calculated Djouadi:1991tk ; Dawson:1990zj
$$\mathrm{\Delta }=\frac{11g^2}{16\pi ^2}.$$
(3)
The full NLO result will require the evaluation of the virtual corrections to the Higgs + 4 parton processes, which are the subject of this paper, the calculation of the tree graph rates from the Higgs \+ 5 partons amplitudes already given in refs. DelDuca:2004wt ; Badger:2004ty ; Dixon:2004za and the calculation of a set of subtraction terms. The subtractions remove singularities present in the real emission diagrams in the regions of soft and collinear emission. After integration over the momentum of the un observed parton they are added back to the virtual emission diagrams and cancel the singularities in those virtual terms.
We believe this calculation would be a useful addition to the literature for several reasons. First, the effective Lagrangian approach appears to be valid for light Higgs boson mass if the transverse momentum of the associated jets is less than the top quark mass DelDuca:2001eu ; DelDuca:2001fn . Second, this process constitutes a โbackgroundโ to the experimentally interesting vector boson fusion process, Fig 1(a). A complete NLO calculation will improve knowledge of this โbackgroundโ process. In addition, because the vector boson fusion process has a well determined normalization, it is one of the most accurate sources of information about the couplings of the Higgs boson at the LHC Zeppenfeld:2000td . An uncontrolled background from gluon fusion process could compromise that measurement. For a comprehensive review of standard model Higgs physics, see ref. Djouadi:2005gi .
Note that the process calculated in this paper is distinguished from the vector boson process, Fig 1(a), by the presence of colored particles exchanged in the $`t`$-channel. The exchange of color charge generates extra jet activity in the central region, allowing discrimination against this process by a jet veto. Although the efficacy of such a veto will finally have to be determined by experiment, it will still be interesting to see how this works at the parton level with a full NLO calculation <sup>1</sup><sup>1</sup>1To a limited extent this has been looked at in ref. DelDuca:2004wt . However in a tree graph calculation one cannot look at the effects of finite jet size or of the central jet veto..
In the large top quark mass limit, virtual corrections have been considered in the effective theory by previous authors. Loop corrections to the process $`Hgg`$ are considered at one-loop level in ref. Dawson:1990zj and at two loop level in refs. Harlander:2000mg ; Anastasiou:2002yz . The results for the processes $`Hggg`$ and $`Hq\overline{q}g`$ are given in refs. Schmidt:1997wr ; Ravindran:2002dc . In the following we shall describe results for the virtual corrections to the processes
$`A)H`$ $``$ $`q\overline{q}q^{}\overline{q}^{},`$ (4)
$`B)H`$ $``$ $`q\overline{q}q\overline{q},`$ (5)
$`C)H`$ $``$ $`q\overline{q}gg,`$ (6)
$`D)H`$ $``$ $`gggg,`$ (7)
using the effective theory, Eq. (1).
## II Lowest order process
### II.1 $`Hq\overline{q}q^{}\overline{q}^{}`$
We first perform the calculation of the matrix element for the process involving two distinct flavors of massless quarks, $`q`$ and $`q^{}`$, process $`A`$,
$$Hq(k_1)+\overline{q}(k_2)+q^{}(k_3)+\overline{q}^{}(k_4).$$
(8)
At Born level, only the diagram in Fig. 2(a) contributes. The color expansion of the amplitude can be written as
$$M_0^A(k_1,k_2,k_3,k_4)=\left[\delta _{i_4}^{i_1}\delta _{i_2}^{i_3}\frac{1}{N_c}\delta _{i_2}^{i_1}\delta _{i_4}^{i_3}\right]a^{(0)}(1,2,3,4),$$
(9)
where $`i_j`$ denotes the color index of the $`j`$th quark and we have introduced the notation
$$a^{(0)}(1,2,3,4)a^{(0)}(k_1,h_1;k_2,h_2;k_3,h_3;k_4,h_4),$$
(10)
where $`k_i`$ and $`h_i`$ denote the momentum and the helicity of quark $`i`$. The result for the squared matrix element summed over the spins and colors of the final state quarks and antiquarks is then
$`A_0(k_1,k_2,k_3,k_4)`$ $``$ $`{\displaystyle |M_0^A(k_1,k_2,k_3,k_4)|^2}`$
$`=`$ $`g^4A^2V\left[{\displaystyle \frac{(s_{13}s_{24}s_{23}s_{14})^2+s_{12}^2s_{34}^2}{s_{34}^2s_{12}^2}}+{\displaystyle \frac{(s_{13}s_{24})^2+(s_{14}s_{23})^2}{2s_{34}s_{12}}}\right].`$
The number of colors, $`N_c`$, enters as $`V=N_c^21`$, so, for the case of SU(3), we have that $`V=8`$. The Lorentz invariants are defined as $`s_{ij}(k_i+k_j)^2=2k_ik_j`$. The momentum of the Higgs can be eliminated in terms of the four massless momenta, $`p_H=k_1k_2k_3k_4`$, so that
$$M_H^2=s_{12}+s_{13}+s_{14}+s_{23}+s_{24}+s_{34}.$$
(12)
### II.2 $`Hq\overline{q}q\overline{q}`$
In the case of massless quarks of identical flavor, process $`B`$,
$$Hq(k_1)+\overline{q}(k_2)+q(k_3)+\overline{q}(k_4),$$
(13)
the Born amplitude squared is determined by the two diagrams, shown in Fig. 2(a) and (b), which differ by the exchange of the final state anti-quarks. The color expansion of the amplitude can be written as
$`M_0^B(k_1,k_2,k_3,k_4)`$ $`=`$ $`\left[\delta _{i_4}^{i_1}\delta _{i_2}^{i_3}{\displaystyle \frac{1}{N_c}}\delta _{i_2}^{i_1}\delta _{i_4}^{i_3}\right]a^{(0)}(1,2,3,4)\left[\delta _{i_2}^{i_1}\delta _{i_4}^{i_3}{\displaystyle \frac{1}{N_c}}\delta _{i_4}^{i_1}\delta _{i_2}^{i_3}\right]a^{(0)}(1,4,3,2)`$ (14)
$`=`$ $`M_0^A(k_1,k_2,k_3,k_4)M_0^A(k_1,k_4,k_3,k_2).`$
The result for the matrix element squared, summed over the spins and colors of the final state quarks and antiquarks is given by
$`B_0(k_1,k_2,k_3,k_4)`$ $``$ $`{\displaystyle |M_0^B(k_1,k_2,k_3,k_4)|^2}`$ (15)
$`=`$ $`A_0(k_1,k_2,k_3,k_4)+A_0(k_1,k_4,k_3,k_2)+B_0^{}(k_1,k_2,k_3,k_4),`$
where the interference term is defined as
$$B_0^{}(k_1,k_2,k_3,k_4)2\mathrm{Re}\left[M_0^A(k_1,k_2,k_3,k_4)^{}M_0^A(k_1,k_4,k_3,k_2)\right].$$
(16)
The result for the lowest order interference term, $`B_0^{}`$, is given by,
$`B_0^{}(k_1,k_2,k_3,k_4)`$ $`=`$ $`g^4A^2C_f\{[(s_{13}s_{24})^2(s_{12}s_{34}+s_{14}s_{23}s_{13}s_{24})`$ (17)
$``$ $`2(s_{13}s_{24}+s_{14}s_{23}s_{12}s_{34})(s_{12}s_{34}+s_{13}s_{24}s_{14}s_{23})]\}`$
$`\times `$ $`{\displaystyle \frac{1}{s_{12}s_{14}s_{23}s_{34}}},`$
with $`C_f=(N_c^21)/(2N_c)=4/3`$.
### II.3 $`Hq\overline{q}gg`$
We now turn to process $`C`$,
$$Hq(k_1)+\overline{q}(k_2)+g(k_3)+g(k_4).$$
(18)
At lowest order the amplitude is given by,
$`M_0^C`$ $`=`$ $`(T^{a_3}T^{a_4})_{i_1i_2}c_1^{(0)}(1,2,3,4)+(T^{a_4}T^{a_3})_{i_1i_2}c_2^{(0)}(1,2,3,4),`$ (19)
where $`a_3,a_4`$ are the color indices of the gluons and $`i_1,i_2`$ are the color indices of the quarks. As before we have introduced notation of the form
$$c_i^{(0)}(1,2,3,4)c_i^{(0)}(k_1,h_1;k_2,h_2;k_3,\epsilon _3;k_4,\epsilon _4),$$
(20)
where $`\epsilon _i`$ is the polarization vector of gluon $`i`$ and $`c_2^{(0)}(1,2,3,4)=c_1^{(0)}(1,2,4,3)`$. Explicit forms for the three independent helicity amplitudes can be found, for example in refs. Kauffman:1997ix ; DelDuca:2004wt . The former reference also contains explicit results for the amplitude squared.
### II.4 $`Hgggg`$
Lastly we consider the matrix element for the process $`D`$,
$$Hg(k_1)+g(k_2)+g(k_3)+g(k_4).$$
(21)
At lowest order the four gluon matrix element has the structure
$`M_0^D`$ $`=`$ $`{\displaystyle \underset{\sigma S_4/Z_4}{}}\text{tr}(T^{a_{\sigma (1)}}T^{a_{\sigma (2)}}T^{a_{\sigma (3)}}T^{a_{\sigma (4)}})d_1^{(0)}(\sigma (1),\sigma (2),\sigma (3),\sigma (4)),`$ (22)
where the sum runs over the six non-cyclic permutations and we have introduced the notation
$$d_i(1,2,3,4)d_i(k_1,\epsilon _1;k_2,\epsilon _2;k_3,\epsilon _3;k_4,\epsilon _4).$$
(23)
The partial amplitudes satisfy the relations Mangano:1987xk ; Berends:1988me
$`d_i^{(0)}(1,2,3,4)=d_i^{(0)}(4,1,2,3)`$ $`\text{cyclicity},`$ (24)
$`d_i^{(0)}(1,2,3,4)=d_i^{(0)}(4,3,2,1)`$ $`\text{reflection},`$ (25)
$`d_i^{(0)}(1,2,3,4)+d_i^{(0)}(2,1,3,4)+d_i^{(0)}(2,3,1,4)=0`$ $`\text{dual Ward identity},`$ (26)
so that at Born level for fixed helicities there are only two independent amplitudes. Explicit expressions for the helicity amplitudes can be found for example in refs. Kauffman:1997ix ; DelDuca:2004wt . The former reference also contains explicit results for the amplitude squared. Eqs. (24) and (25) continue to be valid beyond leading order Bern:1990ux .
## III Higher order processes
In order to control the divergences which will occur at higher order we will continue the dimensionality of space-time, $`D=42ฯต`$. Within the context of dimensional regularization there remain choices of the dimensionality of internal and external gluons which are needed to completely specify the scheme. The most commonly adopted choices are the conventional dimensional regularization, (CDR), the โt Hooft-Veltman scheme, (HV) 'tHooft:1972fi , and the four-dimensional helicity scheme, (FDH) Bern:1991aq ; Bern:2002zk . In the CDR scheme one uniformly continues all momenta and polarization vectors to $`D`$ dimensions. The HV scheme differs in the treatment of the external states, which remain four-dimensional. Finally in the FDH scheme all states are four-dimensional, and only the internal loop momenta are continued to $`D`$ dimensions.
Since we are interested in numerical evaluation, it is preferable to consider the external quarks and gluons in four dimensions, with two physical helicity states. We choose to work in the โt Hooft-Veltman scheme. The relationship of the CDR, HV and FDH regularization schemes has been presented in refs. Kunszt:1993sd ; Catani:1996pk . It is therefore straightforward to translate our results to another scheme. The details of the translation between the HV and FDH schemes are provided in Section IV.
### III.1 Distinct quarks
At next-to-leading order in the perturbative expansion, 30 virtual diagrams contribute to the amplitude given in Eq. (8). At one-loop level the amplitude can be decomposed into two independent color structures,
$$M_1^A(k_1,k_2,k_3,k_4)=\left[\delta _{i_4}^{i_1}\delta _{i_2}^{i_3}\frac{1}{N_c}\delta _{i_2}^{i_1}\delta _{i_4}^{i_3}\right]a_1^{(1)}(1,2,3,4)+\delta _{i_2}^{i_1}\delta _{i_4}^{i_3}a_2^{(1)}(1,2,3,4).$$
(27)
The color sub-amplitude $`a_2^{(1)}`$ does not contribute at next-to-leading order because the interference with the color structure of the Born amplitude vanishes.
Before renormalization we find for the squared matrix element, summed over spin and colors of the final state
$`A_1(k_1,k_2,k_3,k_4)`$ $``$ $`{\displaystyle \left(|M_0^A+M_1^A|^2|M_1^A|^2\right)}`$ (28)
$`=`$ $`A_0(k_1,k_2,k_3,k_4)\left[1+{\displaystyle \frac{g^2}{8\pi ^2}}Y^A(k_1,k_2,k_3,k_4)\right]`$
$`+`$ $`A^2{\displaystyle \frac{V}{2}}{\displaystyle \frac{g^6}{8\pi ^2}}[X^A(k_1,k_2,k_3,k_4)+X^A(k_3,k_4,k_1,k_2)`$
$`+`$ $`X^A(k_2,k_1,k_4,k_3)+X^A(k_4,k_3,k_2,k_1)]+O(ฯต).`$
All ultraviolet and infrared singularities are in the functions $`Y(k_1,k_2,k_3,k_4)`$ given by
$`Y^A(k_1,k_2,k_3,k_4)=N_c{\displaystyle \frac{c_\mathrm{\Gamma }\mu ^{2ฯต}}{ฯต^2}}\left[(s_{14})^ฯต+(s_{23})^ฯต\right]`$ (29)
$`+`$ $`{\displaystyle \frac{1}{N_c}}{\displaystyle \frac{c_\mathrm{\Gamma }\mu ^{2ฯต}}{ฯต^2}}[(s_{12})^ฯต+(s_{34})^ฯต2(s_{13})^ฯต+2(s_{14})^ฯต+2(s_{23})^ฯต2(s_{24})^ฯต)]`$
$``$ $`{\displaystyle \frac{c_\mathrm{\Gamma }\mu ^{2ฯต}}{ฯต}}\left[3C_fb_0\right]\left[(s_{12})^ฯต+(s_{34})^ฯต\right]{\displaystyle \frac{20}{9}}n_f+{\displaystyle \frac{152}{9}}N_c16C_f`$
$`+`$ $`{\displaystyle \frac{1}{N_c}}\left[Ls_1^{2\mathrm{m}e}(s_{134},s_{234};s_{34},M_H^2)+Ls_1^{2\mathrm{m}e}(s_{123},s_{124};s_{12},M_H^2)\right]`$
$``$ $`{\displaystyle \frac{2}{N_c}}\left[Ls_1^{2\mathrm{m}e}(s_{123},s_{134};s_{13},M_H^2)+Ls_1^{2\mathrm{m}e}(s_{124},s_{234};s_{24},M_H^2)\right]`$
$``$ $`(N_c{\displaystyle \frac{2}{N_c}})\left[Ls_1^{2\mathrm{m}e}(s_{124},s_{134};s_{14},M_H^2)+Ls_1^{2\mathrm{m}e}(s_{123},s_{234};s_{23},M_H^2)\right],`$
where
$$c_\mathrm{\Gamma }(4\pi )^ฯต\frac{\mathrm{\Gamma }(1+ฯต)\mathrm{\Gamma }^2(1ฯต)}{\mathrm{\Gamma }(12ฯต)}=\frac{(4\pi )^ฯต}{\mathrm{\Gamma }(1ฯต)}+O(ฯต^3),$$
(30)
and
$$b_0=\left(\frac{11N_c}{3}\frac{2n_f}{3}\right).$$
(31)
As usual $`n_f`$ is the number of light flavors and $`\mu `$ is the scale introduced to keep the coupling constant dimensionless in $`D`$ dimensions.
The finite function $`X^A(k_1,k_2,k_3,k_4)`$ is given by
$`X^A(k_1,k_2,k_3,k_4)`$ $`=`$ $`Ls_1(s_{12},s_{13};s_{123}){\displaystyle \frac{2}{N_c}}f_1(k_2,k_1,k_3,k_4)`$
$`+`$ $`Ls_1(s_{12},s_{23};s_{123})(N_c{\displaystyle \frac{2}{N_c}})f_1(k_1,k_2,k_3,k_4)`$
$`+`$ $`({\displaystyle \frac{1}{N_c}}+N_c)\mathrm{L}_1({\displaystyle \frac{s_{123}}{s_{12}}})f_2(k_1,k_2,k_3,k_4)+N_c\mathrm{L}_0({\displaystyle \frac{s_{123}}{s_{12}}})f_3(k_1,k_2,k_3,k_4)`$
$`+`$ $`{\displaystyle \frac{1}{N_c}}\mathrm{L}_0({\displaystyle \frac{s_{123}}{s_{12}}})f_4(k_1,k_2,k_3,k_4)+({\displaystyle \frac{N_c}{2}}+{\displaystyle \frac{1}{N_c}})\mathrm{L}_0({\displaystyle \frac{s_{124}}{s_{14}}})f_5(k_1,k_2,k_3,k_4)`$
$``$ $`{\displaystyle \frac{1}{N_c}}\mathrm{L}_0({\displaystyle \frac{s_{123}}{s_{13}}})f_5(k_1,k_2,k_4,k_3)+N_c\mathrm{ln}({\displaystyle \frac{s_{123}}{s_{12}}})f_6(k_1,k_2,k_3,k_4)`$
$`+`$ $`N_c\mathrm{ln}({\displaystyle \frac{s_{123}}{s_{23}}})f_7(k_1,k_2,k_3,k_4)+N_c\mathrm{ln}({\displaystyle \frac{s_{12}}{s_{14}}})f_8(k_1,k_2,k_3,k_4)`$
$`+`$ $`{\displaystyle \frac{1}{N_c}}\mathrm{ln}({\displaystyle \frac{s_{123}}{s_{12}}})f_9(k_1,k_2,k_3,k_4)+{\displaystyle \frac{1}{N_c}}\mathrm{ln}({\displaystyle \frac{s_{123}}{s_{13}}})f_{10}(k_1,k_2,k_3,k_4)`$
$``$ $`{\displaystyle \frac{1}{N_c}}\mathrm{ln}({\displaystyle \frac{s_{123}}{s_{23}}})f_{10}(k_2,k_1,k_3,k_4)+(N_c+{\displaystyle \frac{1}{N_c}})f_{12}(k_1,k_2,k_3,k_4).`$
$`+`$ $`{\displaystyle \frac{1}{2N_c}}\left(\mathrm{ln}({\displaystyle \frac{s_{12}}{s_{13}}})+\mathrm{ln}({\displaystyle \frac{s_{12}}{s_{14}}})\right)\left(f_{11}(k_1,k_2,k_3,k_4)f_{11}(k_2,k_1,k_3,k_4)\right).`$
The special functions coming from the loop integrals, $`L_0,L_1,Ls_1`$ and $`Ls_1^{2\mathrm{m}e}`$ are given in Appendix A. The explicit expression for the kinematic functions $`f_i`$ are given in Appendix B. We note that the line-reversal symmetry ($`12`$ and $`34`$) and the renaming property ($`13`$ and $`24`$) are manifest in Eq. (28).
The UV divergences are removed in the $`\overline{\text{MS}}`$-scheme by adding a counterterm $`A_{\mathrm{ct}}`$ given by
$$A_{\mathrm{ct}}(k_1,k_2,k_3,k_4)=2\frac{c_\mathrm{\Gamma }}{ฯต}b_0\frac{g^2}{16\pi ^2}A_0(k_1,k_2,k_3,k_4).$$
(33)
Additionally, there is a finite contribution, $`A_{\mathrm{fin}}`$, coming from the effective Lagrangian, Eq. (1), which is
$$A_{\mathrm{fin}}(k_1,k_2,k_3,k_4)=2\mathrm{\Delta }A_0(k_1,k_2,k_3,k_4),$$
(34)
where $`\mathrm{\Delta }`$ is given in Eq. (3).
### III.2 Identical quarks
In the case of identical quarks, 60 diagrams contribute the next-to-leading order process, Eq. (13). Before renormalization we find for the squared amplitude, summed over colors and spins,
$`B_1(k_1,k_2,k_3,k_4)`$ $``$ $`{\displaystyle \left(|M_0^B+M_1^B|^2|M_1^B|^2\right)}`$ (35)
$`=`$ $`A_1(k_1,k_2,k_3,k_4)+A_1(k_1,k_4,k_3,k_2)+B_1^{}(k_1,k_2,k_3,k_4),`$ (36)
with $`A_1`$ given in (28). The result for the interference term can be written as,
$`B_1^{}(k_1,k_2,k_3,k_4)`$ $`=`$ $`B_0^{}(k_1,k_2,k_3,k_4)\left[1+{\displaystyle \frac{g^2}{8\pi ^2}}Y^B(k_1,k_2,k_3,k_4)\right]`$ (37)
$`+`$ $`A^2V{\displaystyle \frac{g^6}{8\pi ^2}}[X^B(k_1,k_2,k_3,k_4)+X^B(k_3,k_2,k_1,k_4)+X^B(k_1,k_4,k_3,k_2)`$
$`+`$ $`X^B(k_3,k_4,k_1,k_2)+X^B(k_4,k_3,k_2,k_1)+X^B(k_2,k_3,k_4,k_1)`$
$`+`$ $`X^B(k_4,k_1,k_2,k_3)+X^B(k_2,k_1,k_4,k_3)]+O(ฯต),`$
where the function $`Y^B`$ contains all divergent terms
$`Y^B(k_1,k_2,k_3,k_4)={\displaystyle \frac{c_\mathrm{\Gamma }N_c\mu ^{2ฯต}}{ฯต^2}}((s_{24})^ฯต+(s_{13})^ฯต)`$ (38)
$`+`$ $`{\displaystyle \frac{c_\mathrm{\Gamma }\mu ^{2ฯต}}{N_cฯต^2}}\left[(s_{12})^ฯต+(s_{34})^ฯต+(s_{14})^ฯต+(s_{23})^ฯต(s_{24})^ฯต(s_{13})^ฯต\right]`$
$`+`$ $`{\displaystyle \frac{c_\mathrm{\Gamma }\mu ^{2ฯต}}{4ฯต}}\left(6C_f+2b_0\right)\left[(s_{12})^ฯต+(s_{14})^ฯต+(s_{23})^ฯต+(s_{34})^ฯต\right]{\displaystyle \frac{20n_f}{9}}+{\displaystyle \frac{80N_c}{9}}+{\displaystyle \frac{8}{N_c}}`$
$`+`$ $`{\displaystyle \frac{1}{N_c}}[Ls_1^{2\mathrm{m}e}(s_{134},s_{234};s_{34},M_H^2)+Ls_1^{2\mathrm{m}e}(s_{123},s_{234};s_{23},M_H^2)`$
$`+`$ $`Ls_1^{2\mathrm{m}e}(s_{124},s_{134};s_{14},M_H^2)+Ls_1^{2\mathrm{m}e}(s_{123},s_{124};s_{12},M_H^2)]`$
$``$ $`(N_c+{\displaystyle \frac{1}{N_c}})\left[Ls_1^{2\mathrm{m}e}(s_{123},s_{134};s_{13},M_H^2)+Ls_1^{2\mathrm{m}e}(s_{124},s_{234};s_{24},M_H^2)\right].`$
The finite function $`X^B`$ is given by
$`X^B(k_1,k_2,k_3,k_4)=Ls_1(s_{12},s_{24};s_{124})g_1(k_1,k_2,k_3,k_4)\left(1+{\displaystyle \frac{1}{N_c^2}}\right)`$ (39)
$``$ $`Ls_1(s_{12},s_{23};s_{123})g_2(k_1,k_2,k_3,k_4){\displaystyle \frac{1}{N_c^2}}`$
$`+`$ $`\mathrm{L}_1\left({\displaystyle \frac{s_{123}}{s_{12}}}\right)g_3(k_1,k_2,k_3,k_4)\left(1+{\displaystyle \frac{1}{N_c^2}}\right)`$
$`+`$ $`\mathrm{L}_0\left({\displaystyle \frac{s_{123}}{s_{12}}}\right)g_4(k_1,k_2,k_3,k_4)+\mathrm{L}_0\left({\displaystyle \frac{s_{123}}{s_{12}}}\right)g_5(k_1,k_2,k_3,k_4){\displaystyle \frac{1}{N_c^2}}`$
$`+`$ $`\mathrm{ln}\left({\displaystyle \frac{s_{123}}{s_{12}}}\right)g_6(k_1,k_2,k_3,k_4)+\mathrm{ln}\left({\displaystyle \frac{s_{123}}{s_{12}}}\right)g_7(k_1,k_2,k_3,k_4){\displaystyle \frac{1}{N_c^2}}`$
$`+`$ $`g_8(k_1,k_2,k_3,k_4)\left(1+{\displaystyle \frac{1}{N_c^2}}\right),`$
where the functions $`g_i`$ are given in Appendix C. We note that the result in Eq. (35) is symmetric under the exchange of $`(13)`$ or $`(24)`$.
The counterterm renormalizing the ultraviolet divergences in the case of identical quarks reads
$$B_{\mathrm{ct}}(k_1,k_2,k_3,k_4)=2\frac{c_\mathrm{\Gamma }}{ฯต}b_0\frac{g^2}{16\pi ^2}B_0(k_1,k_2,k_3,k_4),$$
(40)
while finite contribution coming from the effective Lagrangian is
$$B_{\mathrm{fin}}(k_1,k_2,k_3,k_4)=2\mathrm{\Delta }B_0(k_1,k_2,k_3,k_4).$$
(41)
### III.3 $`Hq\overline{q}gg`$
At one loop the full amplitude is calculated from 191 Feynman diagrams which can be decomposed into the three color-ordered sub-amplitudes,
$$M_1^C=(T^{a_3}T^{a_4})_{i_1i_2}c_1^{(1)}(1,2,3,4)+(T^{a_4}T^{a_3})_{i_1i_2}c_2^{(1)}(1,2,3,4)+\delta ^{a_3a_4}\delta _{i_1i_2}c_3^{(1)}(1,2,3,4).$$
(42)
Bose symmetry requires that $`c_2^{(1)}(1,2,3,4)=c_1^{(1)}(1,2,4,3)`$.
The divergent parts of these one-loop amplitudes are given by
$`c_1^{(1)}(1,2,3,4)`$ $`c_\mathrm{\Gamma }{\displaystyle \frac{g^2\mu ^{2ฯต}}{16\pi ^2}}[{\displaystyle \frac{N_c}{ฯต^2}}((s_{24})^ฯต+(s_{13})^ฯต+(s_{34})^ฯต)+{\displaystyle \frac{1}{N_cฯต^2}}(s_{12})^ฯต`$ (43)
$``$ $`{\displaystyle \frac{3C_f}{ฯต}}+{\displaystyle \frac{b_0}{ฯต}}]c_1^{(0)}(1,2,3,4)`$
$`c_3^{(1)}(1,2,3,4)`$ $`c_\mathrm{\Gamma }{\displaystyle \frac{g^2\mu ^{2ฯต}}{16\pi ^2}}[{\displaystyle \frac{1}{2ฯต^2}}c_1^{(0)}(1,2,3,4)((s_{14})^ฯต+(s_{23})^ฯต(s_{12})^ฯต(s_{34})^ฯต)`$ (44)
$`+`$ $`{\displaystyle \frac{1}{2ฯต^2}}c_2^{(0)}(1,2,3,4)((s_{13})^ฯต+(s_{24})^ฯต(s_{12})^ฯต(s_{34})^ฯต)].`$
The interference between the Born and the NLO amplitude is given by
$`2\text{Re}(M_1^CM_0^C)`$ $`=`$ $`{\displaystyle \frac{VN_c}{2}}\text{Re}[c_1^{(1)}c_1^{(0)}+c_2^{(1)}c_2^{(0)}]`$ (45)
$``$ $`{\displaystyle \frac{V}{2N_c}}\text{Re}[(c_1^{(1)}+c_2^{(1)})(c_1^{(0)}+c_2^{(0)})^{}]+V\text{Re}[c_3^{(1)}(c_1^{(0)}+c_2^{(0)})^{}],`$
with $`c_ic_i(1,2,3,4)`$. Counterterms, analogous to those in Eqs. (33, 34) need to be included to obtain the full renormalized result.
Numerical results, which are given in the following section, were generated using an extension of the method suggested in ref. Giele:2004iy . Analytic expressions for the Feynman graphs are generated using Qgraf Nogueira:1991ex and Form Vermaseren:2000nd . The scalar and tensor integrals appearing in the amplitudes are reduced numerically using the Davydychev reduction for the tensor integrals Davydychev:1991va and a recursive procedure similar to the one proposed in Giele:2004iy to reduce all scalar integrals to a small number of analytically known basis integrals. These are then evaluated numerically as a Laurent series in the $`ฯต`$ parameter<sup>2</sup><sup>2</sup>2The numerical Laurent expansion technique was first used in ref. vanHameren:2005ed . In a more general analytic context it was used by many authors before.. The key point of this method is that a record is kept of all previously computed integrals, so that each scalar integral is computed only once. The result of our procedure is a numerical expression for the scalar and tensor integrals component by component each of which has a Laurent expansion in $`ฯต`$. This method will be described in detail in a later paper egz2 . Numerical or semi-numerical methods have also been described in refs. vanHameren:2005ed ; Binoth:2005ff ; Binoth:2002xh ; Ferroglia:2002mz ; Andonov:2004hi ; Belanger:2003sd ; deDoncker:2004bf ; Nagy:2003qn ; delAguila:2004nf .
### III.4 $`Hgggg`$
At NLO the amplitude for process Eq. (21) requires the calculation of 739 Feynman diagrams, which can expanded in nine color sub-amplitudes
$`M_1^D`$ $`=`$ $`{\displaystyle \underset{\sigma S_4/Z_4}{}}\text{tr}(T^{a_{\sigma (1)}}T^{a_{\sigma (2)}}T^{a_{\sigma (3)}}T^{a_{\sigma (4)}})d_1^{(1)}(\sigma (1),\sigma (2),\sigma (3),\sigma (4))`$ (46)
$`+`$ $`{\displaystyle \frac{1}{N_c}}\text{tr}(T^{a_1}T^{a_2})\text{tr}(T^{a_3}T^{a_4})d_2^{(1)}(1,2,3,4)`$
$`+`$ $`{\displaystyle \frac{1}{N_c}}\text{tr}(T^{a_1}T^{a_3})\text{tr}(T^{a_2}T^{a_4})d_2^{(1)}(1,3,2,4)`$
$`+`$ $`{\displaystyle \frac{1}{N_c}}\text{tr}(T^{a_1}T^{a_4})\text{tr}(T^{a_2}T^{a_3})d_2^{(1)}(1,4,2,3).`$
If we discard diagrams with internal quark loops we have the decoupling identity Bern:1990ux
$$d_2^{(1)}(1,2,3,4)=\underset{\sigma S_4/Z_4}{}d_1^{(1)}(\sigma (1),\sigma (2),\sigma (3),\sigma (4)).$$
(47)
However, at NLO the $`d_2`$ terms in Eq. (46) do not receive contributions from internal fermion loops. This can be easily shown by explicitly examining the diagrams with internal fermionic bubbles, triangles, and boxes. The general expansion can thus be simplified as a consequence of Eq. (47) so that
$`M_1^D`$ $`=`$ $`{\displaystyle \underset{\sigma S_4/Z_4}{}}\text{tr}(T^{a_{\sigma (1)}}T^{a_{\sigma (2)}}T^{a_{\sigma (3)}}T^{a_{\sigma (4)}})d_1^{(1)}(\sigma (1),\sigma (2),\sigma (3),\sigma (4))`$
$`+`$ $`{\displaystyle \frac{1}{N_c}}[\text{tr}(T^{a_1}T^{a_2})\text{tr}(T^{a_3}T^{a_4})+\text{tr}(T^{a_1}T^{a_3})\text{tr}(T^{a_2}T^{a_4})`$
$`+`$ $`\text{tr}(T^{a_1}T^{a_4})\text{tr}(T^{a_2}T^{a_3})]d_2^{(1)}(1,2,3,4).`$ (48)
Using Eq. (III.4) it can be shown that the result for the matrix element squared is
$`|M_0^D+M_1^D|^2|M_1^D|^2={\displaystyle \frac{N_c^2(N_c^21)}{16}}{\displaystyle \underset{\sigma S_4/Z_4}{}}\{|d_1^{(0)}(\sigma (1),\sigma (2),\sigma (3),\sigma (4))|^2`$ (49)
$`+`$ $`2\text{Re}\left[d_1^{(0)}(\sigma (1),\sigma (2),\sigma (3),\sigma (4))^{}d_1^{(1)}(\sigma (1),\sigma (2),\sigma (3),\sigma (4))\right]\}.`$
Counterterms, analogous to those in Eqs. (33, 34) need to be included to obtain the full renormalized result.
Numerical results for this matrix element squared were generated using the method described above. The pole structure for the color sub-amplitude $`d_1`$ has the simple form
$$d_1^{(1)}(1,2,3,4)\frac{c_\mathrm{\Gamma }g^2\mu ^{2ฯต}}{16\pi ^2}\left[\frac{N_c}{ฯต^2}\left((s_{12})^ฯต+(s_{23})^ฯต+(s_{34})^ฯต+(s_{14})^ฯต\right)\right]d_1^{(0)}(1,2,3,4).$$
(50)
## IV Numerical results
In this section we present numerical results for the Born amplitude squared and for its interference with the one-loop matrix element for the four processes of interest, $`A,B,C`$ and $`D`$. We use the following arbitrarily chosen, momentum configuration, where a Higgs boson of unit mass decays into four well separated partons, $`(E,p_x,p_y,p_z)`$:
$$\begin{array}{ccccc}p_H=\hfill & (1.00000000000,\hfill & 0.00000000000,\hfill & 0.00000000000,\hfill & 0.00000000000),\hfill \\ k_1=\hfill & (+0.30674037867,\hfill & 0.17738694693,\hfill & 0.01664472021,\hfill & 0.24969277974),\hfill \\ k_2=\hfill & (+0.34445032281,\hfill & +0.14635282800,\hfill & 0.10707762397,\hfill & +0.29285022975),\hfill \\ k_3=\hfill & (+0.22091667641,\hfill & +0.08911915938,\hfill & +0.19733901856,\hfill & +0.04380941793),\hfill \\ k_4=\hfill & (+0.12789262211,\hfill & 0.05808504045,\hfill & 0.07361667438,\hfill & 0.08696686795).\hfill \end{array}$$
(51)
For each process, {$`A,B,C,D`$}, we introduce the quantities
$`X_B`$ $`=`$ $`{\displaystyle \frac{1}{g^4A^2}}X_0(k_1,k_2,k_3,k_4),`$
$`X_V`$ $`=`$ $`{\displaystyle \frac{8\pi ^2}{g^6A^2}}\left[X_1(k_1,k_2,k_3,k_4)X_0(k_1,k_2,k_3,k_4)\right],\text{with }X=A,B,C,D,`$ (52)
which are independent of the value of the coupling constant. Thus $`X_B`$ is the matrix element squared evaluated using the Born amplitude. $`X_{V,N}`$ and $`X_{V,A}`$ denote the contributions of the interference between the virtual amplitude and the lowest order, as calculated from the numerical and analytical formulas. The unrenormalized results are given in Table 1 for the scale choice $`\mu =M_H`$ and the momenta of Eq. (51).
The explicit results show that far from exceptional momentum configurations, where divergent inverse Gram determinants are known to spoil the accuracy of the numerical procedure, a relative accuracy of $`๐ช\left(10^{13}\right)`$ can be achieved. For processes $`C`$ and $`D`$, where a full analytical result is not available, we verified that the answer satisfies the Ward identities to a similar relative accuracy. For process $`D`$ we checked numerically that for $`n_f=0`$, the color amplitudes satisfy the decoupling identity, Eq. (47). Close to exceptional momentum configurations, it is still possible to use a numerical approachGiele:2004ub ; egz2 .
We have also checked numerically that our results satisfy the following relationship between the HV and FDH regularization schemes,
$`a_1^{(1)\text{FDH}}(1,2,3,4)a_1^{(1)\text{HV}}(1,2,3,4)`$ $`=`$ $`{\displaystyle \frac{g^2}{16\pi ^2}}\left({\displaystyle \frac{N_c}{3}}{\displaystyle \frac{1}{N_c}}\right)a^{(0)}(1,2,3,4),`$
$`a_2^{(1)\text{FDH}}(1,2,3,4)a_2^{(1)\text{HV}}(1,2,3,4)`$ $`=`$ $`0,`$
$`c_1^{(1)\text{FDH}}(1,2,3,4)c_1^{(1)\text{HV}}(1,2,3,4)`$ $`=`$ $`{\displaystyle \frac{g^2}{16\pi ^2}}\left({\displaystyle \frac{N_c}{6}}{\displaystyle \frac{1}{2N_c}}\right)c_1^{(0)}(1,2,3,4),`$
$`c_3^{(1)\text{FDH}}(1,2,3,4)c_3^{(1)\text{HV}}(1,2,3,4)`$ $`=`$ $`0,`$
$`d_1^{(1)\text{FDH}}(1,2,3,4)d_1^{(1)\text{HV}}(1,2,3,4)`$ $`=`$ $`0.`$ (53)
Applying the finite renormalization which compensates for the difference between the ultraviolet regularization in the two schemes Kunszt:1993sd , we recover the expected difference between the two schemes due to the differing infrared regularization,
$`a_1^{(1)\text{FDH}}(1,2,3,4)a_1^{(1)\text{HV}}(1,2,3,4)`$ $`=`$ $`{\displaystyle \frac{g^2}{4\pi ^2}}\stackrel{~}{\gamma }_qa^{(0)}(1,2,3,4),`$
$`c_1^{(1)\text{FDH}}(1,2,3,4)c_1^{(1)\text{HV}}(1,2,3,4)`$ $`=`$ $`{\displaystyle \frac{g^2}{8\pi ^2}}(\stackrel{~}{\gamma }_q+\stackrel{~}{\gamma }_g)c_1^{(0)}(1,2,3,4),`$
$`d_1^{(1)\text{FDH}}(1,2,3,4)d_1^{(1)\text{HV}}(1,2,3,4)`$ $`=`$ $`{\displaystyle \frac{g^2}{4\pi ^2}}\stackrel{~}{\gamma }_gd_1^{(0)}(1,2,3,4),`$ (54)
where
$$\stackrel{~}{\gamma }_q\frac{C_f}{2}\text{and}\stackrel{~}{\gamma }_g\frac{N_c}{6}.$$
(55)
The other two relations in Eq. (53) are unchanged.
## V Outlook
In this paper we presented results obtained using a general, semi-numerical calculation of one-loop corrections. In order to establish the feasibility of the semi-numerical method, we computed all the one-loop corrections to Higgs plus four parton processes using an effective Lagrangian. We presented explicit results for a specific, non-exceptional phase space point. For practical applications of this method, one has to be able to treat exceptional momentum configurations also. The method of this paper can be extended to the treat those regions. A detailed description of the algorithm is presented in a separate work egz2 .
The results presented in this paper generate two separate lines of research. The first is clearly the completion of the calculation of the Higgs boson plus two jet process at next-to-leading order. As indicated in the text all of the needed elements are now in place.
The second development is the extension of these methods to calculate other one-loop processes which currently lie beyond the range of analytic calculation. Examples of processes of current experimental interest are diboson plus one jet ($`V_1,V_2,j`$), tri-boson production ($`V_1,V_2,V_3`$) and vector boson plus heavy quark pairs ($`VQ\overline{Q}`$).
### Acknowledgments
We are happy to acknowledge useful discussions with W.A. Bardeen, E.W.N. Glover and U. Haisch. We thank Carola Berger and Lance Dixon for pointing out typos in the analytical expressions of four quark amplitudes, which have been corrected in the present version of the paper.
## Appendix A Integral Functions Appearing in Amplitudes
The integral functions appearing in the virtual corrections are presented in this appendix. Following closely the notation of ref. Bern:1997sc we define
$$\mathrm{L}_0(r)=\frac{\mathrm{ln}(r)}{1r},\mathrm{L}_1(r)=\frac{\mathrm{L}_0(r)+1}{1r}.$$
(56)
The above functions have the property that they are finite as their denominators vanish. Furthermore we define
$`Ls_1(s,t;m^2)`$ $`=`$ $`\mathrm{Li}_2(1{\displaystyle \frac{s}{m^2}})+\mathrm{Li}_2(1{\displaystyle \frac{t}{m^2}})+\mathrm{ln}{\displaystyle \frac{s}{m^2}}\mathrm{ln}{\displaystyle \frac{t}{m^2}}{\displaystyle \frac{\pi ^2}{6}},`$ (57)
where the dilogarithm is defined as usual as
$$\mathrm{Li}_2(x)=_0^x๐z\frac{\mathrm{ln}(1z)}{z}.$$
(58)
The function $`Ls_1`$ is simply related to the scalar box integral with one external mass evaluated in six space-time dimensions, where it is infrared- and ultraviolet-finite.
The โeasyโ six-dimensional box function with two non-adjacent external masses, $`m_1,m_3`$, is related to the function $`Ls_1^{2\mathrm{m}e}`$
$`Ls_1^{2\mathrm{m}e}(s,t;m_1^2,m_3^2)`$ $`=\mathrm{Li}_2\left(1{\displaystyle \frac{m_1^2}{s}}\right)\mathrm{Li}_2\left(1{\displaystyle \frac{m_1^2}{t}}\right)\mathrm{Li}_2\left(1{\displaystyle \frac{m_3^2}{s}}\right)\mathrm{Li}_2\left(1{\displaystyle \frac{m_3^2}{t}}\right)`$ (59)
$`+\mathrm{Li}_2\left(1{\displaystyle \frac{m_1^2m_3^2}{st}}\right){\displaystyle \frac{1}{2}}\mathrm{ln}^2\left({\displaystyle \frac{s}{t}}\right).`$ (60)
This function has the property that it vanishes as $`s+tm_1^2m_3^20`$. The analytic continuation of these integrals is obtained adding a small positive imaginary part to each invariant, $`s_{ij}s_{ij}+i\epsilon `$.
## Appendix B Functions for distinct quarks
The kinematic functions for the virtual corrections to $`Hq\overline{q}q^{}\overline{q}^{}`$ appearing in Eq. (III.1) are given below:
$`f_1(k_1,k_2,k_3,k_4)`$ $`=`$ $`{\displaystyle \frac{s_{12}s_{34}}{2s_{13}^2}}+{\displaystyle \frac{3s_{13}s_{24}s_{23}^2+s_{14}s_{23}s_{14}^2s_{13}^2}{s_{12}s_{34}}}{\displaystyle \frac{s_{14}^2s_{23}^2}{2s_{12}s_{13}^2s_{34}}}`$ (61)
$``$ $`{\displaystyle \frac{s_{24}^2}{2s_{12}s_{34}}}2{\displaystyle \frac{(s_{13}s_{24}s_{14}s_{23})^2}{s_{12}^2s_{34}^2}}+{\displaystyle \frac{s_{24}}{s_{13}}}+{\displaystyle \frac{s_{14}s_{23}}{s_{13}^2}}2`$
$`f_2(k_1,k_2,k_3,k_4)`$ $`=`$ $`{\displaystyle \frac{s_{12}s_{34}(s_{12}s_{34}+s_{23}(s_{24}+2s_{23}s_{14}))+s_{23}^2(s_{24}+s_{14})^2}{2s_{12}^3s_{34}}}`$ (62)
$`f_3(k_1,k_2,k_3,k_4)`$ $`=`$ $`{\displaystyle \frac{s_{34}}{2s_{23}}}+s_{23}(s_{24}+s_{14}){\displaystyle \frac{s_{24}+4s_{23}+3s_{14}}{2s_{12}^2s_{34}}}+{\displaystyle \frac{3s_{24}+4s_{23}}{2s_{12}}}`$ (63)
$`f_4(k_1,k_2,k_3,k_4)`$ $`=`$ $`2{\displaystyle \frac{s_{34}}{s_{23}}}s_{23}(s_{24}+s_{14}){\displaystyle \frac{s_{24}+2s_{23}+5s_{14}}{2s_{12}^2s_{34}}}`$ (64)
$``$ $`{\displaystyle \frac{4s_{24}+6s_{23}3s_{14}}{2s_{12}}}`$
$`f_5(k_1,k_2,k_3,k_4)`$ $`=`$ $`{\displaystyle \frac{s_{13}}{s_{14}}}{\displaystyle \frac{2s_{23}}{s_{24}}}+{\displaystyle \frac{s_{24}}{s_{14}}}{\displaystyle \frac{s_{23}}{s_{34}}}{\displaystyle \frac{s_{14}s_{23}^2}{s_{24}^2s_{34}}}+{\displaystyle \frac{s_{13}s_{23}}{s_{24}s_{34}}}+{\displaystyle \frac{s_{13}s_{24}}{s_{14}s_{34}}}{\displaystyle \frac{s_{34}}{s_{14}}}`$ (65)
$`f_6(k_1,k_2,k_3,k_4)`$ $`=`$ $`{\displaystyle \frac{s_{12}s_{34}}{2s_{13}s_{23}}}+{\displaystyle \frac{4s_{23}s_{24}+2s_{14}s_{24}3s_{13}s_{24}+3s_{14}s_{23}}{2s_{12}s_{34}}}`$ (66)
$`+`$ $`{\displaystyle \frac{s_{14}^2s_{23}}{2s_{12}s_{13}s_{34}}}{\displaystyle \frac{s_{14}}{s_{13}}}+{\displaystyle \frac{1}{2}}`$
$`f_7(k_1,k_2,k_3,k_4)`$ $`=`$ $`{\displaystyle \frac{s_{14}^2s_{23}^2+s_{13}s_{34}^2s_{12}s_{13}s_{14}s_{23}s_{24}}{2s_{13}^2s_{34}s_{12}}}`$ (67)
$`f_8(k_1,k_2,k_3,k_4)`$ $`=`$ $`{\displaystyle \frac{s_{14}s_{23}s_{13}s_{24}}{2s_{12}s_{34}}}`$ (68)
$`f_9(k_1,k_2,k_3,k_4)`$ $`=`$ $`{\displaystyle \frac{s_{13}s_{24}^2}{s_{12}s_{23}s_{34}}}{\displaystyle \frac{s_{12}s_{34}}{s_{13}s_{23}}}{\displaystyle \frac{s_{14}^2s_{23}}{s_{12}s_{13}s_{34}}}2{\displaystyle \frac{s_{24}}{s_{23}}}+2{\displaystyle \frac{s_{14}}{s_{13}}}1`$ (69)
$``$ $`{\displaystyle \frac{2s_{24}^2+2s_{23}s_{24}+5s_{14}s_{24}5s_{13}s_{24}+5s_{14}s_{23}}{2s_{12}s_{34}}}`$
$`f_{10}(k_1,k_2,k_3,k_4)`$ $`=`$ $`{\displaystyle \frac{s_{12}s_{23}s_{34}^2+s_{13}^2s_{24}^2s_{13}s_{14}s_{23}s_{24}}{s_{12}s_{23}^2s_{34}}}`$ (71)
$`f_{11}(k_1,k_2,k_3,k_4)`$ $`=`$ $`{\displaystyle \frac{2s_{13}s_{24}}{s_{12}s_{34}}}`$ (72)
$`f_{12}(k_1,k_2,k_3,k_4)`$ $`=`$ $`{\displaystyle \frac{s_{13}(s_{13}(s_{14}s_{24})+2s_{14}s_{23})}{2s_{12}^2s_{34}}}+{\displaystyle \frac{s_{14}}{2s_{12}}}.`$ (73)
## Appendix C Functions for identical quarks
The kinematic functions for the virtual corrections to $`Hq\overline{q}q\overline{q}`$ appearing in Eq. (35) are given below:
$`g_1(k_1,k_2,k_3,k_4)`$ $`=`$ $`1{\displaystyle \frac{s_{13}s_{24}\left(s_{13}^2+s_{24}^2\right)}{4s_{12}s_{14}s_{23}s_{34}}}+{\displaystyle \frac{s_{13}^2+2s_{14}s_{23}2s_{13}s_{24}+s_{24}^2}{4s_{12}s_{34}}}`$ (74)
$`+`$ $`{\displaystyle \frac{s_{13}^22s_{13}s_{24}+s_{24}^2+2s_{12}s_{34}}{4s_{14}s_{23}}}`$
$`g_2(k_1,k_2,k_3,k_4)`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \frac{s_{24}}{4s_{13}}}{\displaystyle \frac{s_{14}^2s_{23}^2}{8s_{12}s_{13}^2s_{34}}}+{\displaystyle \frac{3s_{14}s_{23}s_{24}}{8s_{12}s_{13}s_{34}}}+{\displaystyle \frac{s_{13}s_{24}\left(s_{13}^2+s_{24}^2\right)}{4s_{12}s_{14}s_{23}s_{34}}}`$ (75)
$``$ $`{\displaystyle \frac{2s_{13}^2+3s_{14}s_{23}3s_{13}s_{24}+4s_{24}^2}{8s_{12}s_{34}}}+{\displaystyle \frac{3s_{12}s_{24}s_{34}}{8s_{13}s_{14}s_{23}}}{\displaystyle \frac{s_{12}^2s_{34}^2}{8s_{13}^2s_{14}s_{23}}}`$
$`+`$ $`{\displaystyle \frac{s_{14}s_{23}+s_{12}s_{34}}{8s_{13}^2}}{\displaystyle \frac{2s_{13}^23s_{13}s_{24}+4s_{24}^2+3s_{12}s_{34}}{8s_{14}s_{23}}}`$
$`g_3(k_1,k_2,k_3,k_4)`$ $`=`$ $`{\displaystyle \frac{s_{23}^2s_{24}}{8s_{12}^2s_{14}}}+{\displaystyle \frac{s_{23}\left(s_{14}+s_{23}+s_{24}\right)}{8s_{12}^2}}{\displaystyle \frac{s_{34}}{4s_{12}}}+{\displaystyle \frac{\left(s_{23}+s_{24}\right)s_{34}}{8s_{12}s_{14}}}+{\displaystyle \frac{s_{34}^2}{8s_{14}s_{23}}}`$ (76)
$`g_4(k_1,k_2,k_3,k_4)`$ $`=`$ $`{\displaystyle \frac{6s_{23}3s_{24}}{8s_{12}}}+{\displaystyle \frac{s_{24}\left(s_{23}+4s_{24}\right)}{8s_{12}s_{14}}}{\displaystyle \frac{s_{34}}{4s_{14}}}+{\displaystyle \frac{5s_{24}s_{34}}{8s_{14}s_{23}}}`$ (77)
$`g_5(k_1,k_2,k_3,k_4)`$ $`=`$ $`{\displaystyle \frac{s_{14}}{4s_{12}}}+{\displaystyle \frac{s_{24}}{8s_{12}}}{\displaystyle \frac{s_{24}^2}{4s_{12}s_{14}}}+{\displaystyle \frac{3s_{23}\left(2s_{14}+s_{24}\right)}{8s_{12}s_{14}}}{\displaystyle \frac{s_{23}\left(s_{14}+s_{24}\right)^2}{4s_{12}^2s_{34}}}`$ (78)
$`+`$ $`{\displaystyle \frac{s_{34}}{4s_{14}}}+{\displaystyle \frac{s_{34}}{4s_{23}}}{\displaystyle \frac{3s_{24}s_{34}}{8s_{14}s_{23}}}{\displaystyle \frac{s_{12}s_{34}^2}{4s_{14}s_{23}^2}}`$
$`g_6(k_1,k_2,k_3,k_4)`$ $`=`$ $`{\displaystyle \frac{5}{8}}{\displaystyle \frac{s_{24}}{4s_{14}}}+{\displaystyle \frac{3s_{13}s_{24}+4s_{24}^23s_{12}s_{34}}{8s_{14}s_{23}}}`$ (79)
$`g_7(k_1,k_2,k_3,k_4)`$ $`=`$ $`{\displaystyle \frac{7}{8}}+{\displaystyle \frac{s_{14}}{4s_{13}}}+{\displaystyle \frac{s_{24}}{4s_{14}}}{\displaystyle \frac{s_{14}^2s_{23}}{4s_{12}s_{13}s_{34}}}{\displaystyle \frac{s_{24}^2}{4s_{12}s_{34}}}+{\displaystyle \frac{s_{12}s_{34}}{4s_{13}s_{23}}}{\displaystyle \frac{s_{12}^2s_{34}^2}{4s_{13}s_{14}s_{23}^2}}`$ (80)
$``$ $`{\displaystyle \frac{s_{13}s_{24}+2s_{24}^2s_{12}s_{34}}{8s_{14}s_{23}}}`$
$`g_8(k_1,k_2,k_3,k_4)`$ $`=`$ $`{\displaystyle \frac{s_{12}}{32s_{14}}}+{\displaystyle \frac{s_{14}}{32s_{12}}}+{\displaystyle \frac{s_{12}}{32s_{23}}}+{\displaystyle \frac{s_{13}\left(s_{14}2s_{24}\right)}{64s_{12}s_{23}}}{\displaystyle \frac{s_{13}s_{24}}{32s_{12}s_{14}}}`$ (81)
$`+`$ $`{\displaystyle \frac{s_{13}s_{24}\left(s_{13}+s_{24}\right)}{64s_{12}s_{14}s_{23}}}+{\displaystyle \frac{s_{23}\left(2s_{14}+s_{24}\right)}{64s_{12}s_{14}}}+{\displaystyle \frac{s_{14}+s_{23}}{32s_{34}}}+{\displaystyle \frac{s_{13}\left(s_{12}+s_{23}2s_{24}\right)}{64s_{14}s_{34}}}`$
$`+`$ $`{\displaystyle \frac{\left(s_{12}2s_{13}+s_{14}\right)s_{24}}{64s_{23}s_{34}}}+{\displaystyle \frac{s_{13}s_{24}\left(s_{13}+s_{24}\right)}{64s_{12}s_{14}s_{34}}}+{\displaystyle \frac{s_{13}s_{24}\left(s_{13}+s_{24}\right)}{64s_{12}s_{23}s_{34}}}`$
$`+`$ $`{\displaystyle \frac{s_{13}s_{24}\left(s_{13}+s_{24}\right)}{64s_{14}s_{23}s_{34}}}+{\displaystyle \frac{\left(s_{13}s_{14}+2s_{12}\left(s_{14}+s_{23}\right)+s_{23}s_{24}\right)s_{34}}{64s_{12}s_{14}s_{23}}}.`$
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# Renormalization of NN Interaction with Chiral Two Pion Exchange Potential. Central Phases and the Deuteron.
## I Introduction
The possibility suggested by Weinberg Weinberg (1990) and pioneered by Ray, Ordoรฑez and Van Kolck Ordonez and van Kolck (1992); Van Kolck (1993); Ordonez et al. (1996) of making model independent predictions for NN scattering using Effective Field Theory (EFT) methods and, more specifically, Chiral Perturbation Theory (ChPT) has triggered a lot of activity in recent years (for a review see e.g. Ref. Bedaque and van Kolck (2002)). In addition to the previous works, most subsequent calculations dealing with the specific consequences of ChPT have been focused in making predictions for NN scattering phase-shifts and deuteron properties based in the genuine Two Pion Exchange (TPE) chiral potentials Rijken and Stoks (1996); Kaiser et al. (1997, 1998); Epelbaum et al. (1998, 2000); Rentmeester et al. (1999); Friar (1999); Richardson (1999); Kaiser (2000a, b, 2002, 2001a, 2001b); Entem and Machleidt (2002a, b); Rentmeester et al. (2003); Epelbaum et al. (2004a, b); Entem and Machleidt (2003a); Higa and Robilotta (2003); Higa et al. (2004); Higa (2004) although some incipient work has also recently been started implementing Three Pion Exchange effects Kaiser (2000a, b, 2001b); Entem and Machleidt (2003b); Epelbaum et al. (2005). In a given partial wave (coupled) channel with good total angular momentum the reduced NN potential ($`U(r)=MV(r)`$) in configuration space can schematically be written in local and energy independent form Kaiser et al. (1997); Friar (1999); Rentmeester et al. (1999) (see Refs. Ordonez and van Kolck (1992); Van Kolck (1993); Ordonez et al. (1996) for an energy dependent representation) for any distances larger than a finite short distance radial regulator, $`r_c`$,
$`U(r)`$ $`=`$ $`{\displaystyle \frac{Mm^3}{f^2}}W_{\mathrm{LO}}(mr,g)+{\displaystyle \frac{Mm^5}{f^4}}W_{\mathrm{NLO}}(mr,g,\overline{d})`$ (1)
$`+`$ $`{\displaystyle \frac{m^6}{f^4}}W_{\mathrm{NNLO}}(mr,g,\overline{c}_1,\overline{c}_3,\overline{c}_4)+\mathrm{}`$
where $`W(x)`$ are known dimensionless functions which are everywhere finite except for the origin where they exhibit power law divergencies, which demand the use of some regularization. In writting the previous expression we have disregarded distributional contact terms (deltas and derivatives of deltas) which strengths are scheme dependent but do not contribute for $`rr_c>0`$. The potential is completely specified by the pion mass, $`m`$, the pion weak decay constant, $`f`$, the nucleon mass $`M`$, the axial coupling constant $`g`$, the Goldberger-Treimann discrepancy $`\overline{d}_{18}`$ and three additional low energy constants $`\overline{c}_1=c_1M`$, $`\overline{c}_3=c_3M`$ and $`\overline{c}_4=c_4M`$ which can directly be deduced from the analysis of low energy $`\pi N`$ scattering within ChPT Fettes et al. (1998); Buettiker and Meissner (2000); Gomez Nicola et al. (2000); Nicola et al. (2004). Given this information one can then solve the single or coupled channel Schrรถdinger equation imposing a regularity condition of the wave function at the origin for each separate channel. In this paper we will work under the assumption that the long range pieces of the potential should be iterated to all orders, but some perturbative analysis will also be done. Our motivation is to describe long range correlations between observables in the NN problem in a model independent way. As we will see below there is still some additional physical information required in the form of either counter-terms or short distance boundary conditions to make the problem well posed if one indeed wants to go to remove the regularization. They depend exclusively on the short distance behaviour of the potential through phases of the wave function. The number of phases depends crucially on the repulsive or attractive character of the potential. In this sense the power counting for the short distance interactions cannot be regarded as independent on the power counting of the singular chiral potentials. Non-perturbatively this materializes, after renormalization, in non-integer power counting for physical observables. This imposes severe limitations on the admissible structure of counterterms and the corresponding renormalization conditions of the quantum mechanical problem. Right away we hast to emphasize that our approach is not the conventional EFT one of allowing all possible short distance counterterms allowed by the symmetry, and to a certain extent our viewpoints are admittedly heterodox within the conventional EFT framework. However, based on the physical requirement of having small wave functions in the short range unknown region the basic and orthodox quantum mechanical requirements of completeness and orthoghonality of states are deduced, providing a justification for the additional restrictions. We remind here the series of works by Phillips and Cohen Phillips and Cohen (1997, 2000) (see also Re. Scaldeferri et al. (1997)), where restrictions on zero range interactions were deduced for non-singular potentials based on the Wigner causality conditions. Here we extend their results also to the singular NN interactions of Eq. (1).
The theorem underlying the EFT developments is that if chiral symmetry is spontaneously broken down in QCD, the true NN potential at long distances is embedded in the parameter envelope of the general chiral NN potential, Eq. (1), and the chiral expansion provides a reliable hierarchy at those long distances. The hope is that compatible and perhaps accurate determinations of both $`\pi N`$ and $`NN`$ low energy data, bound states and resonances can be achieved with the same sets of parameters. The problem is that, in order to make truly model independent predictions, short distance ambiguities should be under control and their size smaller than the experimental data uncertainties used as input of the calculation. Only then can the renormalization program be carried out satisfactorily as it was done in the OPE case Frederico et al. (1999); Beane et al. (2002); Pavon Valderrama and Ruiz Arriola (2004a, b, 2005a); Nogga et al. (2005), although, as recognized by Nogga, Timmermans and Van Kolck, this may be done at the expense of modifying the power counting Nogga et al. (2005) of the counterterms in favor of renormalizability (see also Ref. Pavon Valderrama and Ruiz Arriola (2005a) for a complementary formulation in terms of boundary conditions). The present work analyses this problem extending our previous OPE renormalized calculations to NNLO TPE and its implication in the values of the chiral constants.
The determination of the chiral constants $`c_1`$, $`c_3`$ and $`c_4`$ (in units of $`\mathrm{GeV}^1`$ from now on) from $`\pi N`$ scattering has been undertaken in several works and shows significant systematic discrepancies depending on the details of the analysis. In Heavy Baryon ChPT for low energy $`\pi N`$ scattering Fettes et al. (1998) the values $`c_1=1.23\pm 0.16`$, $`c_3=5.94\pm 0.09`$ and $`c_4=3.47\pm 0.05`$ were deduced with a sigma term of $`\sigma (0)=70\mathrm{M}\mathrm{e}\mathrm{V}`$. In Ref. Buettiker and Meissner (2000) the analysis of low energy $`\pi N`$ scattering inside the Mandelstam triangle yields $`c_1=0.81\pm 0.15`$, $`c_3=4.69\pm 1.34`$ and $`c_4=3.40\pm 0.04`$ with, however, a bit too low sigma term $`\sigma (0)=40\mathrm{M}\mathrm{e}\mathrm{V}`$ as compared to ChPT. Unitarization methods reproducing the phase-shifts Gomez Nicola et al. (2000); Nicola et al. (2004) from threshold to the $`\mathrm{\Delta }`$ resonance region conclude $`c_1=0.43\pm 0.04`$, $`c_3=3.10\pm 0.05`$ and $`c_4=1.51\pm 0.04`$.
The values of the chiral constants $`c_1`$, $`c_3`$ and $`c_4`$ also depend on regularization details of the $`NN`$ chiral interaction. The $`\pi N`$ values from Ref. Buettiker and Meissner (2000) were taken in the nucleon-nucleon NNLO calculation of Ref. Epelbaum et al. (2000) with sharp and gaussian cut-offs $`\mathrm{\Lambda }=0.60.8\mathrm{GeV}`$ in momentum space, and momentum dependent counter-terms were supplemented and determined from a fit to the NN data base based on the Partial Wave Analysis (PWA) of Ref. Stoks et al. (1993, 1994). Likewise, Ref. Entem and Machleidt (2002a) constructs a NNLO chiral potential where channel dependent gaussian momentum space cut-offs in the range $`\mathrm{\Lambda }=0.40.5\mathrm{GeV}`$ were used to fit the NN database Arndt et al. (1994). The N<sup>3</sup>LO extension of this work Entem and Machleidt (2003b) uses only one common cut-off and fixing $`c_1=0.81`$ produces $`c_3=3.20`$ and $`c_4=5.40`$. In Ref. Rentmeester et al. (1999) the NNLO calculation was done in configuration space with a short distance cut-off at $`r=1.4\mathrm{fm}`$ where an energy and channel dependent boundary condition was imposed and the fixed value $`c_1=0.76\pm 0.07`$ was used to make a PWA to pp data yielding $`c_3=5.08\pm 0.24`$ and $`c_4=4.70\pm 0.70`$. An update of this calculation also including np data Rentmeester et al. (2003) generates $`c_3=4.78\pm 0.10`$ and $`c_4=3.96\pm 0.22`$. The calculations of Ref. Epelbaum et al. (2004a, b) improve the cut-off dependence of the potential in momentum space by using spectral regularization and take again gaussian cut-offs and fix $`c_1=0.81`$ yielding, after fitting the counter-terms to the NN PWA Stoks et al. (1993, 1994), the values $`c_3=3.40`$ and $`c_4=3.40`$. The extension of this work to N<sup>3</sup>LO has been done in Ref. Epelbaum et al. (2005) keeping the same values for $`c_3`$ and $`c_4`$ and readjusting the counter-terms.
In a renormalized theory results should be insensitive to the auxiliary regularization method if the regulator is removed at the end. If a fit to the database proves successful, then the resulting parameters should be cut-off independent or at least the systematic uncertainty induced by the regularization should be smaller than the statistical errors induced by experimental data. Otherwise, the cut-off becomes a physical parameter. The first indication that finite cut-offs effects are sizeable in present calculations has to do with the variety of values that have been used in the literature for the low energy constants $`c_1`$, $`c_3`$ and $`c_4`$ to adjust NN partial waves and deuteron properties Rentmeester et al. (1999); Epelbaum et al. (2000); Epelbaum et al. (2004a, b); Entem and Machleidt (2003b); Rentmeester et al. (2003) (see also the comment in Ref.Entem and Machleidt (2003a)). Obviously, we do not expect the values of the cโs to agree exactly, but the discrepancies should be at the level of the difference in the approximation <sup>1</sup><sup>1</sup>1For instance, pion loops at NLO modify the contribution to $`c_3`$ by $`3g_A^2m^2/(64\pi f^2)0.4/\mathrm{GeV}`$. This contribution must be taken into account when comparing numbers between Rentmeester et al. (1999, 2003), Buettiker and Meissner (2000) and the present approach. Only the extractions using the N3LO $`NN`$ potential in Ref. Entem and Machleidt (2003b) are made at the same order as those from Buettiker and Meissner (2000). . Since the data base is the same but the regularization schemes are different, one unavoidably suspects that these determinations of the LECS may perhaps be regularization and hence cut-off dependent.
To get a proper perspective on the issue of renormalization let us consider the size of the contributions of the chiral potential in configuration space at different distances. For instance, at $`r=1.4\mathrm{fm}`$ in the $`{}_{}{}^{1}S_{0}^{}`$ channel each order in the expansion is about an order of magnitude smaller than the preceding one. At short distances, however, the situation is exactly the opposite, higher orders dominate over the lower orders. In the previous example of the $`{}_{}{}^{1}S_{0}^{}`$ channel, LO and NLO become comparable at $`r0.9\mathrm{fm}`$, and NLO and NNLO become comparable at distances which value $`r0.10.4\mathrm{fm}`$ depends strongly on the particular choice of low energy constants $`c_3`$ and $`c_4`$. Actually, a general feature of the chiral NN potentials at NNLO has to do with their short distance behavior; they develop an attractive Van der Waals singularity $`UMC_6/r^6`$ similar to the one found for neutral atomic systems. In such a situation, the standard regularity condition at the origin only specifies the wave function uniquely if the potential is repulsive, but some additional information is required if the potential is attractive Case (1950) (for a comprehensive review in the one channel case see e.g. Ref. Frank et al. (1971).). Within the EFT framework the problem has been revisited in Ref. Beane et al. (2001). The net result is that the regularity condition at the origin tames the singularity Pavon Valderrama and Ruiz Arriola (2005a) and, in fact, more singular potentials become less important at low energies.
In this work we reanalyze the NN chiral potential including TPE potential at NNLO. We carry out the analysis entirely in coordinate space following the ideas developed in our previous work Pavon Valderrama and Ruiz Arriola (2005a) for the OPE potential. In configuration space the (renormalized) potential is finite except at the origin, a point which should carefully be handled, requiring a delicate numerical limiting procedure. For a singular potential in coordinate space, the corresponding potential in momentum space is not finite unless a short distance cut-off or a subtraction procedure at the level of the potential is implemented, hence modifying the potential everywhere and not just at high energies. This results generally in two cut-offs: one for the irreducible two point function and another for the Lippmann-Schwinger iteration Epelbaum et al. (2000); Entem and Machleidt (2002a, b); Epelbaum et al. (2004a, b); Entem and Machleidt (2003b); Epelbaum et al. (2005). The short distance coordinate space cut-off is unique and common both to the potential and the scattering solution.
Unlike previous works on the TPE potential we try to remove the cut-off completely taking the consequences seriously. This does not mean that finite cut-off calculations are necessarily incorrect or not entitled to describe all or some part of the data, but there are also good reasons for removing the cut-off and looking at the physical consequences. In the first place the limit exists in strict mathematical sense under well defined conditions, as the analysis below shows. This is a non-trivial fact, because calculations done in momentum space can only address this question numerically by adding counterterms suggested by an a priori power counting on the short distance potential. As shown in Ref. Nogga et al. (2005) this does not always work, and calculations may require some trial and error. Secondly, this is the only way we know how to get rid of short distance ambiguities, and thus to make truly model independent calculations. Third, the study of peripheral waves has proven to be successful by using perturbative renormalized amplitudes corresponding to irreducible TPE and iterated OPE where the cut-off has been removed in the intermediate state Kaiser et al. (1997); Entem and Machleidt (2002b). Peripheral waves mainly probe large distances in the Born approximation but they also see some of the short distance interaction due to re-scattering effects. Fourth, the advantage of renormalization is that one should obtain the same results provided one uses as input the same physical information, regardless whether the calculation is done in coordinate or momentum space, and also regardless on the particular regularization. Finally, a reliable estimate on the errors and convergence rate of the chiral expansion can be done, without any spurious cut-off contamination. In principle, the higher order in the chiral expansion the better, provided there is perfect errorless data to fit the increasing number of low energy constants appearing at any order. However, the chiral expansion may reach a limited predictive power because of finite experimental accuracy in the low energy constants used as input. The output inherits a propagated error which may eventually become larger than the experimental uncertainty <sup>2</sup><sup>2</sup>2This issue has been illustrated in Refs. Nieves and Ruiz Arriola (2000); Colangelo et al. (2000, 2001) for the case of $`\pi \pi `$ scattering at two loops, and will become clear in NN scattering below.. Finite cut-off uncertainties are not a substitute for propagating input experimental errors to the predictions of the theory, and can be regarded at best as a lower bound on systematic errors. In this paper we regard this possible cut-off dependence as purely numerical inaccuracies of the calculation, and not as a measure of the uncertainty in the predictions of the theory, so we make any effort to minimize these cut-off induced systematic errors.
In the process of eliminating the cut-off we find some surprises, and effects not explored up to now become manifest. Even for low energy scattering parameters and deuteron properties, where the description should be more reliable and robust we find systematic discrepancies in our calculation with values quoted in the literature and which we conclusively identify as finite cut-off effects. This might provide a natural explanation why calculations with different cut-off methods fitting the NN phase shifts Stoks et al. (1993, 1994); Arndt et al. (1994) obtain different results for the chiral constants $`c_1`$, $`c_3`$ and $`c_4`$ or why different values of the constants yield good fits to the data. According to our study, for the lowest phases the reason can partly be related to the dominance of short distance Van der Waals singularities for a system with unnaturally large scattering lengths or a weakly bound state as it is the case for the $`{}_{}{}^{1}S_{0}^{}`$ and $`{}_{}{}^{3}S_{1}^{}^3D_1`$ channels. In some cases, they are as large as a $`30\%`$ effect like in the effective range of the triplet $`{}_{}{}^{3}S_{1}^{}`$ channel. The size of the effect depends on the value of the low energy $`\pi N`$ constants $`c_1`$, $`c_3`$ and $`c_4`$. Given the significant sensitivity of low energy NN properties and deuteron properties on these low energy $`\pi N`$ constants we try to make a fit to some low energy properties which uncertainties are reliably known and where we expect the chiral theory to be most reliable. At this point we depart from the standard large scale fits to all phase shifts or partial wave analysis where the low energy threshold parameters are determined a posteriori. The assignment of statistical errors on the fitting parameters $`c_1`$, $`c_3`$ and $`c_4`$ is often not addressed (see however Refs. Rentmeester et al. (1999, 2003)) because the NN data bases used to fit the phase shifts Stoks et al. (1993, 1994); Arndt et al. (1994) are treated as errorless. We also try to improve on this point within our framework.
The paper is organized as follows. In Sect. II we discuss the basic assumption of the smallness of the wave function in the short range unknown region and its consequences. We also analyze the constraints based on causality and analyticity of the $`S`$matrix. In Sect. III we introduce the classification of boundary conditions which will be used along the paper to effectively renormalize the amplitudes both in the one-channel as well as the coupled channel case. We will also review the orthogonality constraints for singular potentials already used in our previous work Pavon Valderrama and Ruiz Arriola (2005a) for the OPE potential. Sect. IV deals with the description of the singlet $`{}_{}{}^{1}S_{0}^{}`$ channel. From the superposition principle of boundary conditions we show how a universal form of a low energy theorem for the threshold parameters as well as for the phase-shift arises. In Sect. V we discuss the interesting triplet $`{}_{}{}^{3}S_{1}^{}^3D_1`$ channel both for the deuteron bound state as well as the corresponding scattering states, where full use of the orthogonality constraints as well as the superposition principle of boundary conditions generates interesting analytical relations connecting deuteron and scattering properties. In Sect. VI a careful discussion of errors for our cut-off independent results is carried out. Also, a determination of the chiral constants based on low energy data and deuteron properties is made. In Sect. VII we present a simplified study on the significance of the chiral Van der Waals forces and the striking similarities with the full calculations for the s-waves. In Sect. VIII we show some puzzling results for the NLO calculation in the deuteron channel. We also comment the relation to finite cut-off calculations and the conflict between Weinberg counting and non-perturbative renormalization at NLO. We also outline possible solutions to this problem. In Sect. IX we analyze our results on the light of long distance perturbation theory, reinforcing the usefulness of non-perturbative renormalization due to an undesirable proliferation of counterterms. Finally, in Sect. X we summarize our conclusions.
For numerical calculations we take $`f_\pi =92.4\mathrm{MeV}`$, $`m=138.03\mathrm{MeV}`$, $`M=M_pM_n/(M_p+M_n)=938.918\mathrm{MeV}`$, $`g_{\pi NN}=13.083`$ in the OPE piece to account for the Goldberger-Treimann discrepancy according to the Nijmegen phase shift analysis NN scattering de Swart et al. (1997) and $`g_A=1.26`$ in the TPE piece of the potential. The values of the coefficients $`c_1`$, $`c_3`$ and $`c_4`$ used along this paper are listed in Table 1 for completeness. The potentials in configuration space used in this paper are exactly those provided in Ref. Kaiser et al. (1997); Friar (1999); Rentmeester et al. (1999) but disregarding relativistic corrections, $`M/E1`$.
## II Short distance insensitivity conditions and Renormalization
In this section we ellaborate on the essential role played by standard quantum mechanical orthogonality and completeness properties of the wave functions in the rest of this paper. As we have already mentioned, our approach is unconventional from an EFT perspective, and at present it is unclear whether such properties have an EFT justification. At the same time one should say that many self-denominated EFT calculations do indeed normalize deuteron wave functions to unity and use energy independent regulators from which orthogonality relations follow automatically.
### II.1 The inner and the outer regions
Similar to EFT, our basic assumption is that low energy physics should not depend on short distance fine details. This rather general principle can be made into a precise quantitative statement in practice for a quantum mechanical system. For the sake of clarity let us consider the singlet $`{}_{}{}^{1}S_{0}^{}`$ channel for positive energies. If we assume a short distance regulator $`r_c`$, above which our long distance (local) potential acts, the reduced Schrรถdinger equation in the outer region reads
$`u_{k,\mathrm{L}}^{\prime \prime }(r)+U_\mathrm{L}(r)u_{k,\mathrm{L}}(r)=k^2u_{k,\mathrm{L}}(r),r>r_c`$ (2)
where the label L stands for long. Asymptotically behaves as
$`u_{k,\mathrm{L}}(r)A\mathrm{sin}(kr+\delta (k,r_c))`$ (3)
where $`A`$ is an arbitrary normalization constant and the dependence on the short distance regulator $`r_c`$ has been explicitly highlighted. In the inner region, the dynamics is unknown but we also expect it to be irrelevant provided $`kr_c1`$, i.e. if we assume the corresponding wavelength to be larger than the short distance scale. The potential can be deduced from perturbation theory in the full amplitude,
$`U(\stackrel{}{x})=C_0\delta (\stackrel{}{x})+C_2\{^2,\delta (\stackrel{}{x})\}+\mathrm{}+U_\mathrm{L}(x)`$ (4)
where $`U_L(x)`$ corresponds to the expansion in Eq. (1Kaiser et al. (1997); Friar (1999); Rentmeester et al. (1999). The distributional contact terms are regularization scheme dependent and correspond to polynomial terms in momentum space. Obviously, they do not contribute to the region $`r>r_c`$ for $`r_c>0`$. The very nature of such a calculation already implies that $`C_0`$, $`C_2`$, etc. are perturbative corrections to the short range physics, but do not include possible non-perturbative effects. As will become clear below (Sect. III), finiteness of the physical phase shitf in the limit $`r_c0`$, implies a highly non-perturbative reinterpretation of the short range terms, even if the long range pieces are computed perturbatively.
Following Phillips and Cohen (1997); Scaldeferri et al. (1997) it is useful to use a nonlocal and energy independent potential to describe the short distance dynamics
$`u_{k,\mathrm{S}}^{\prime \prime }(r)+{\displaystyle _0^{r_c}}U_\mathrm{S}(r,r^{})u_{k,\mathrm{S}}(r^{})`$ $`=`$ $`k^2u_{k,\mathrm{S}}(r),`$ (5)
$`r<r_c`$
where the label S stands for short. This holds provided we are below any inelastic channel such as $`\pi NN`$. Above such threshold, any open channel should be included explicitly. The nonlocal short distance potential $`U_\mathrm{S}(r,r^{})`$ encodes, in particular, contact terms (deltas and derivatives of deltas) which appear when the long distance potential $`U_\mathrm{L}(r)`$ is computed in perturbation theory. These terms are in fact ambiguous (and hence unphysical) and depend on the regularization scheme used in the perturbative calculation, but are not essential since they do not contribute to the absortive part and hence to the corresponding spectral function Kaiser et al. (1997). The important point is that these ambiguous distributional terms never contribute to the long range part provided $`r_c>0`$ since in this case a compact support for distributional terms is guaranteed. This is a clear advantage of the radial cut-off we are using. In fact, our main motivation for carrying out the analysis in coordinate space is this clean separation between short and long distances. Many regulators, mainly those in momentum space, do not fulfill this condition, since the regulator effectively smears these short distance terms and intrudes somewhat into the long distance region <sup>3</sup><sup>3</sup>3For instance, the widely used gaussian regulator in momentum space Entem and Machleidt (2003b); Epelbaum et al. (2005); Nogga et al. (2005) of a local potential in cordinate space corresponds indeed to a convolution of the original potential smeared over the region of size $`1/\mathrm{\Lambda }`$..
Finally, we need a matching condition connecting the inner and outer regions,
$`{\displaystyle \frac{u_{k,\mathrm{L}}^{}(r_c)}{u_{k,\mathrm{L}}(r_c)}}={\displaystyle \frac{u_{k,\mathrm{S}}^{}(r_c)}{u_{k,\mathrm{S}}(r_c)}}`$ (6)
Viewed from the outer region this relation corresponds to an energy dependent boundary condition at a given short distance cut-off radius, $`r_c`$. Because of elastic unitarity we expect the state to be normalized, so that if we use a box of size $`a`$ as an infrared regulator, we have
$`{\displaystyle _0^{r_c}}u_{k,\mathrm{S}}(r)^2๐r+{\displaystyle _{r_c}^a}u_{k,\mathrm{L}}(r)^2๐r=1`$ (7)
with $`a`$ much larger than the range of the potential. This equation gives a quantitative separation between long distance and known physics and short distance and unknown physics. Obviously, any effective description based on the long range part should fulfill
$`{\displaystyle _0^{r_c}}u_{k,\mathrm{S}}(r)^2๐r1`$ (8)
which corresponds to the requirement of a small wave function in the inner unknown region. This is our basic condition from which most of our results follow. It should be realized that here long and short distances are intertwined through the matching condition, Eq. (6). In particular, an arbitrarily growing function at the origin cannot fulfill this condition even if value of the short distance cut-off is taken.
This kind of pathological situation actually occurrs when dealing with the deuteron channel in the theory with no explicit pions Phillips and Cohen (2000), with OPE potential Pavon Valderrama and Ruiz Arriola (2005a) and TPE potential at NLO in the Weinberg counting (see Sect. VIII and VIII.3 below). As an instructive and enlightening example, we illustrate the situation in Fig. 1 in the deuteron state for the quantity
$`P(r_c)`$ $`=`$ $`{\displaystyle _0^{r_c}}(u(r)^2+w(r)^2)๐r`$ (9)
(for notation see Sect. V) using the LO, NLO and NNLO potentials, Eq. (1), in the outer region, $`r>r_c`$, and matching to a free particle in the inner region $`r<r_c`$ (the precise form of the wave function inside turns out not to be essential) <sup>4</sup><sup>4</sup>4In a momentum space formulation this is somewhat equivalent to cut-off the Lippmann-Schwinger equation above a given value $`\mathrm{\Lambda }`$.. As one sees, there are cases such as a pure short distance theory where the asymptotic $`D/S`$ ratio, $`\eta `$, and the deuteron binding energy are fixed to their experimental value. For $`r_c<1.4\mathrm{fm}`$ one has a miminum probability of about $`20\%`$ in the inner region, and then $`P(r_c)`$ starts to grow. In this case the description cannot be considered effective below that critical value. A similar situation occurs in the NLO TPE case where one has a decreasing inside probability until one reaches $`r_c0.8\mathrm{fm}`$ where it takes its minimum value, about $`7\%`$ and then starts increasing. In contrast, for the LO (OPE) case where the deuteron binding energy is fixed to the experimental one and $`\eta `$ is predicted Pavon Valderrama and Ruiz Arriola (2005a) to have the value $`\eta =0.02633`$ from the regularity condition of the wave function at the origin and also the NNLO (TPE) case where both numbers are fixed to experiment, the probability in the interior region is controlled and generally decreasing. This is a first and transparent illustration that the description based on any preconceived power counting is not necessarily consistent with the fact that short distance ambiguities are under control, since the wave function in the inner and unknown region does not become arbitrarily small as the cut-off is removed. Another possibility is to keep a finite cut-off, and we will analyze this point in more detail in Sect. VIII.3.
As we will discuss in the next Sect. III, tight constraints on the structure of short distance counterterms must be fulfilled if the requirement $`P(r_c)0`$ as $`r_c0^+`$ is imposed.
### II.2 Wigner bounds on the short distance contributions
To proceed further, we derive with respect to the energy Eq. (5) and one immediately gets after some algebra,
$`{\displaystyle \frac{d}{dk^2}}\left[{\displaystyle \frac{u_{k,\mathrm{S}}^{}(r_c)}{u_{k,\mathrm{S}}(r_c)}}\right]={\displaystyle \frac{_0^{r_c}u_{k,\mathrm{S}}(r)^2๐r}{u_{k,\mathrm{S}}(r_c)^2}}0`$ (10)
whenever $`r_c`$ is not a zero of the wave function.
The short range theory can be characterized by an accumulated phase shift $`\delta _S(k,r_c)`$, given by the solution of the truncated short range problem,
$`u_{k,\mathrm{S}}(r)=\mathrm{sin}\left(kr+\delta _S(k,r_c)\right),r>r_c`$ (11)
which fulfills, from Eq. (10),
$`{\displaystyle \frac{d}{dk^2}}\left[k\mathrm{cot}\left(kr_c+\delta _S(k,r_c)\right)\right]0`$ (12)
a condition equivalent to Wignerโs causality condition Phillips and Cohen (1997). Using an effective range expansion for the short distance phase shift
$`k\mathrm{cot}\delta _S={\displaystyle \frac{1}{\alpha _{0,S}}}+{\displaystyle \frac{1}{2}}r_{0,S}k^2+\mathrm{}`$ (13)
we get the Wigner bound for the effective range
$`r_{0,S}2r_c\left[1{\displaystyle \frac{r_c}{\alpha _{0,S}}}+{\displaystyle \frac{r_c^2}{3\alpha _{0,S}}}\right]`$ (14)
where $`\alpha _{0,S}`$ and $`r_{0,S}`$ represent the scattering length and effective range when the short distance potential $`U_S(r,r^{})`$ is switched on from the origin up to the scale $`r_c`$ or equivalently when the long distance potential is switched off from infinity down to $`r_c`$ (see Refs. Pavon Valderrama and Ruiz Arriola (2004a, b) for more details). In a theory where the long distant potential is absent $`U_\mathrm{L}(r)=0`$, i.e., a pure short distance description, we have the obvious result that the short distance threshold parameters coincide with the physical parameters $`\alpha _{0,S}=\alpha _0`$ and $`r_{0,S}=r_0`$. Thus,
$`r_02r_c\left[1{\displaystyle \frac{r_c}{\alpha _0}}+{\displaystyle \frac{r_c^2}{3\alpha _0}}\right],U_\mathrm{L}(r)=0`$ (15)
which implies that $`r_00`$ for $`r_c0`$. With the experimental values one gets the lowest short distance cut-off compatible with causality to be $`r_c=1.4\mathrm{fm}`$. For a given long distance potential we just solve the equations from infinity inwards and look for the point where the Wigner condition is first violated. In Fig. 2 we plot the evolution of the Wigner bound on the effective range for the different approximations to the potential according to the expansion (1) as a function of the short distance cut-off radius. As we see, the lower bound on the radius is pushed towards the origin, and in fact for the NNLO approximation there is no lower bound at all. Thus, only for the NNLO TPE potential can one build the full strength of the experimental effective range without violation of the Wigner condition. We will see more on this in Sect. IV.
So far the discussion has been carried out for a fixed value of $`r_c`$. If we change the short distance radius, $`r_cr_c+\mathrm{\Delta }r_c`$, we can use the matching condition, Eq. (6), to evaluate the change seen from the outer region. This results into a variable phase equation which has been analyzed extensively in our previous works Pavon Valderrama and Ruiz Arriola (2004a, b). If we take the limit $`r_c0`$, we get that the outer wave functions fulfills
$`{\displaystyle \frac{d}{dk^2}}\left[{\displaystyle \frac{u_k^{}(0^+)}{u_k(0^+)}}\right]=0`$ (16)
where the label L has been supressed. Thus, the boundary condition becomes energy independent when the limit $`r_c0^+`$ is taken if the inner wave function becomes arbitrarily small. Note that the limit is taken from above such that $`r_c>0`$ <sup>5</sup><sup>5</sup>5To illustrate this point, let us note that for potentials which do not diverge too strongly at the origin $`r^2U(r)0`$ such as OPE in the $`{}_{}{}^{1}S_{0}^{}`$ channel, there is some irreversibility in the process of integrating from exactly $`r_c=0`$ out (which requires the regular solution, $`u(0)=0`$) or integrating in towards the origin from above $`r0^+`$ (which generally involves the irregular solution, $`u(0)0`$). See the discussion in Ref. Pavon Valderrama and Ruiz Arriola (2004a) and in Sect. III.3. As we have said already, this justifies not considering contact terms in the potential. A direct consequence is that we get also that with the exception of the short distance scattering length $`\alpha _{S,0}(0^+)`$, which can be fixed by some renormalization condition (see Sect. III), the remaining short distance threshold parameters are also zero in this limit,
$`r_{0,S}(0^+)=0v_{2,S}(0^+)=0,\mathrm{}`$ (17)
This energy independence of the boundary condition at the origin insures the orthogonality conditions between different energy states,
$`{\displaystyle _{0^+}^{\mathrm{}}}u_k(r)u_k^{}(r)๐r=\delta (kk^{}).`$ (18)
## III Short distance behavior of CHIRAL POTENTIALS, orthogonality constraints and the number of independent constants
As we have said, chiral NN potentials, Eq. (1), although decay exponentially at large distances, become singular at short distances, where one has
$`U(r)={\displaystyle \frac{MC_n}{r^n}}\left(1+a_1r+a_2r^2+\mathrm{}\right)`$ (19)
To avoid any misconception let us emphasize that the short distance behaviour of a long distance potential should be regarded as a long distance feature, i.e. a long wavelength property, since different long distance potentials yield different short distance behaviours. The short distance properties of chiral potentials have nothing to do with short distance properties of the โtrueโ potential, but renormalization and finiteness requires a very precise behaviour of the wave function when approaching the origin from long distances. In this section we classify the undetermined constants depending on the attractive or repulsive nature of the corresponding potentials in the single channel case and the eigenvalues of the potential matrix in the coupled channel case. In coordinate space and disregarding relativistic corrections the potentials in Eq. (1) are local and energy independent <sup>6</sup><sup>6</sup>6In momentum space and up to NNLO the long distance part of the potential depends on the momentum transfer $`q`$ only and not on the total momentum $`k`$. Essential non-localities, i.e. contributions of the form $`V(q,k)=L(q)k^2`$ with $`L(q)`$ a non-polynomial function, depend weakly on the total momentum and appear first at N<sup>3</sup>LO Kaiser (2000a, b, 2002, 2001a, 2001b); Entem and Machleidt (2002b) due to relativistic $`1/M^2`$ one loop contributions. In coordinate space this weak non-locality corresponds to a modification of the kinetic energy term in the form of a general self-adjoint Sturm-Liouville operator, $`u^{\prime \prime }(r)(p(r)u^{}(r))^{}`$, with a singular $`p(r)`$ function at the origin and exponentially decaying at long distances. The present formalism can in principle be extended to include these features, and will be discussed elsewhere. Nevertheless, according to the results of Sect. VI (see Table 4 ) on the loss of predictive power already at NNLO there is a lack of phenomenological motivation.. An important condition on the short distance behavior of the wave functions are the orthogonality constraints between states of different energy. For a regular energy independent potential these constraints are automatically satisfied, but for singular potentials they generate new relations relevant to the NN interaction. Our approach is not the conventional one of adding short distance counterterms following an a priori power counting on the long distance potential. Rather it is the power counting in the potential what uniquely determines the admissible form of the short distance physics if we want to reach a finite limit when the regulator is removed. This can only be achieved by choosing the regular solution at the origin, i.e. $`u(0)=0`$. In Sect. VIII.3 we will show that irregular solutions generate divergent results after renormalization. Although this may look as a potential drawback of removing the cut-off it provides valuable insight on the form of the potential and the validity of the expansion (see Sect. VIII.4).
### III.1 One channel case
Let us first review the single channel case in a way that results for the coupled channel situation can be easily stated. The reduced Schrรถdinger equation for angular momentum $`l`$ is
$`u^{\prime \prime }+U(r)u+{\displaystyle \frac{l(l+1)}{r^2}}=k^2u.`$ (20)
For a power law singular potential at the origin of the form $`U(r)=MC_n/r^n=\pm (R/r)^n/R^2`$ with $`n>2`$ and $`R`$ the length scale dimension, the de Broglie wavelength is given by $`1/k(r)=1/\sqrt{|U(r)|}`$ and the applicability condition for the WKB approximation reads $`(1/k(r))^{}1`$, so that for distances $`rR(n/2)^{2/(2+n)}`$ one has a semiclassical wave function Case (1950); Frank et al. (1971); Beane et al. (2001). Keeping the leading short distance behavior one gets for attractive and repulsive singular potentials and any angular momentum the following behaviour for the regular solutions
$`u_A(r)`$ $``$ $`C_A\left({\displaystyle \frac{r}{R}}\right)^{n/4}\mathrm{sin}\left[{\displaystyle \frac{2}{n2}}\left({\displaystyle \frac{R}{r}}\right)^{\frac{n}{2}1}+\phi \right],`$
$`\mathrm{for}U_A`$ $``$ $`{\displaystyle \frac{1}{R^2}}\left({\displaystyle \frac{R}{r}}\right)^n`$ (22)
$`u_R(r)`$ $``$ $`C_R\left({\displaystyle \frac{r}{R}}\right)^{n/4}\mathrm{exp}\left[{\displaystyle \frac{2}{n2}}\left({\displaystyle \frac{R}{r}}\right)^{\frac{n}{2}1}\right],`$ (23)
$`\mathrm{for}U_R`$ $``$ $`+{\displaystyle \frac{1}{R^2}}\left({\displaystyle \frac{R}{r}}\right)^n,`$ (24)
respectively. Here $`C_A`$ and $`C_R`$ are normalization constants and $`\phi `$ an arbitrary short distance phase. In the repulsive case we have discarded the irregular solution (a similar exponential with a positive sign) which would not allow to normalize states. For an attractive singular potential there is a short distance unknown parameter. This phase could, in principle, be energy dependent. Chiral potentials are, however, local and energy independent at NNLO at all distances, and become genuinely energy dependent at N<sup>3</sup>LO, due to relativistic $`1/M^2`$ corrections Kaiser (2000a, b, 2002, 2001a, 2001b) <sup>7</sup><sup>7</sup>7The subthreshold energy dependence from the virtual pion production channel $`NNNN\pi `$ which is in principle N<sup>3</sup>LO, disappears since in the heavy baryon limit the threshold $`s_{\pi NN}=(2M+m)^2=4(M^2+k^2)`$ translates into a CM momentum $`k=\sqrt{m(M+m/4)}\mathrm{}`$.. Thus, if we require orthogonality of states with different energy (positive or negative) we get
$`0`$ $`=`$ $`u_k^{}u_pu_pu_k^{}|_0`$ (25)
$`=`$ $`{\displaystyle \frac{1}{R}}\mathrm{sin}\left(\phi (k)\phi (p)\right).`$
Hence, the phase $`\phi `$ is energy independent and could be fixed by matching the solution to the asymptotic large distance region (we assume a short range potential), e.g., by requiring a given value of the scattering length, $`\alpha _l`$, at zero energy. In this way, a new and physical scale appears into the problem which is not specified by the potential. This is equivalent to the well known phenomenon of dimensional transmutation. Another possibility is to fix $`\phi `$ from a given bound state energy, $`E=B`$. The new scale entering the problem is the corresponding wave number, $`\gamma =\sqrt{MB}`$. Note that although neither $`\alpha _l`$ nor $`\gamma `$ can be predicted from a singular potential, the orthogonality constraint does predict a correlation between them through the potential. Likewise, the phase shifts $`\delta _l`$ can be deduced from either $`\alpha _l`$ or $`\gamma `$ by taking the same short distance phase $`\phi `$. In the repulsive case there is no dimensional transmutation since the orthogonality condition follows from regularity at the origin, and the potential fully specifies the wave function. In this case, the scattering length and the spectrum are completely determined from the potential as for standard regular potentials.
### III.2 Coupled channel case
We turn now to the two coupled channel case where the wave functions are denoted by a column vector $`(u,w)`$ (for some particular cases see e.g. Refs. Sprung et al. (1994); Beane et al. (2001); Pavon Valderrama and Ruiz Arriola (2005a)). If we assume that at short distances the reduced potential behaves as
$`UM{\displaystyle \frac{๐_๐ง}{r^n}}`$ (26)
where $`๐_n`$ is a symmetric matrix of Van der Waals coefficients. Diagonalizing the matrix $`๐_๐ง`$ we get
$`๐_๐ง`$ $`=`$ $`\left(\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta \\ \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right)\left(\begin{array}{cc}C_{n,+}& 0\\ 0& C_{n,}\end{array}\right)\left(\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta \\ \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right)`$
(27)
where $`C_{n,\pm }`$ are the corresponding eigenvalues and $`\theta `$ the mixing angle. Thus, at short distances we can decouple the equations to get
$`\left(\begin{array}{c}u\\ w\end{array}\right)`$ $``$ $`\left(\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta \\ \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right)\left(\begin{array}{c}u_+\\ u_{}\end{array}\right)`$ (28)
where $`(u_+,u_{})`$ are regular solutions as in the single channel case. So, in the two channel situation we have three possible cases depending upon the sign of the eigenvalues.
1. Both eigenvalues are negative, i.e., both eigenpotentials are attractive and $`MC_{n,+}=R_+^{n2}`$ and $`MC_{n,}=R_{}^{n2}`$ with $`R_\pm `$ the corresponding scale dimension. In this case the short distance eigensolutions are oscillatory and there are two undetermined short distance phases, $`\phi _+`$ and $`\phi _{}`$. Moreover for two states, $`(u_k,w_k)`$ and $`(u_p,w_p)`$, with different energies we get the orthogonality constraint
$`0`$ $`=`$ $`u_k^{}u_pu_pu_k^{}+w_k^{}w_pw_pw_k^{}|_0`$
$`=`$ $`{\displaystyle \frac{1}{R_+}}\mathrm{sin}\left(\phi _+(k)\phi _+(p)\right)+{\displaystyle \frac{1}{R_{}}}\mathrm{sin}\left(\phi _{}(k)\phi _{}(p)\right).`$
2. One eigenvalue is negative and the other is positive, $`MC_{n,+}=R_+^{n2}`$ and $`MC_{n,}=R_{}^{n2}`$. One short distance eigensolution is a decreasing exponential and the other is oscillatory, so we have one short distance phase $`\phi `$. In this case for two states $`(u_k,w_k)`$ and $`(u_p,w_p)`$ with different energies we get the orthogonality constraint
$`0`$ $`=`$ $`u_k^{}u_pu_pu_k^{}+w_k^{}w_pw_pw_k^{}|_0`$ (30)
$`=`$ $`{\displaystyle \frac{1}{R_+}}\mathrm{sin}\left(\phi _+(k)\phi _+(p)\right).`$
3. Both eigenvalues are positive, $`MC_{n,+}=R_+^{n2}`$ and $`MC_{n,}=R_{}^{n2}`$. Then, both short distance eigensolutions are decreasing exponentials. There are no short distance phases. In this case the orthogonality relations are automatically satisfied.
This simple argument can be easily generalized to any number of coupled channels. The number of undetermined short distance phases corresponds to the number of attractive eigenpotentials at short distances. Orthogonality of the wave functions requires that all these short distance phases fulfill a generalized condition of the form of Eq. (LABEL:eq:orth).
The orthogonality conditions require the determination of the short distance phases, as we did in Ref. Pavon Valderrama and Ruiz Arriola (2005a) for the OPE case. This requires in general an improvement on the short distance behaviour to high orders. An alternative method is to impose the orthogonality constraints either in the single or coupled channel case by integrating in from infinity for a fixed energy, either positive or negative, and then impose the condition at a sufficiently short distance cut-off radius $`r=r_c`$. In the single channel case one would get the condition,
$`{\displaystyle \frac{u_k^{}(r_c)}{u_k(r_c)}}={\displaystyle \frac{u_0^{}(r_c)}{u_0(r_c)}},`$ (31)
if the zero energy state is taken as the reference state. An analogous relation holds for the coupled channel situation, namely
$`0`$ $`=`$ $`u_k(r_c)u_0(r_c)^{}u_k(r_c)^{}u_0(r_c)`$ (32)
$`+`$ $`w_k(r_c)w_0(r_c)^{}w_k(r_c)^{}w_0(r_c).`$
Obviously, in this procedure cut-off independence must be checked. For the TPE chiral potentials analyzed in this paper we find that $`r_c=0.10.2\mathrm{fm}`$ proves a sufficiently small value of the short distance cut-off.
### III.3 Power Counting, counterterms and short distance parameters
As we see, the number of independent parameters is determined from the potential, although their value can be fixed arbitrarily, by some renormalization condition like, e.g., fixing scattering lengths to their physical value. This removes the cut-off in a way that short distances become less and less important. Now, if the potential is regular, i.e., $`r^2|U(r)|0`$, one may choose between the regular and irregular solution <sup>8</sup><sup>8</sup>8Both cases comply to the normalizability condition at the origin, Eq. (8). In the first case the scattering length is predicted while in the second case the scattering length becomes an input of the calculation. In either case the wave function is still normalizable at the origin. Singular potentials at the origin, i.e. fullfiling, $`r^2|U(r)|\mathrm{}`$, do not allow this choice if one insists on normalizability of the wave function at the origin. If the potential is repulsive, the scattering length is fixed while for an attractive potential the scattering length must be an input parameter. Furthermore, orthogonality of different energy solutions requires an energy independence of the boundary condition, so that in all cases the effective range, and higher order threshold parameters cannot be taken as independent parameters, in addition to the scattering lengths.
This can be translated into the language of counterterms quite straightforwardly. In momentum space, fixing $`\alpha _0`$ arbitrarily corresponds to take a constant $`C_0`$ cut-off dependent and energy independent contribution to the potential $`V_0(k^{},k)`$ in the Lippmann-Schwinger equation. Likewise, fixing $`r_0`$ can be mapped as adding a term $`C_2(k^2+k_{}^{}{}_{}{}^{2})`$ to the potential. For higher coupled channel partial waves one fixes the scattering length $`\alpha _{l,l^{}}`$ one has instead terms of the form $`C_{l^{},l}k^l^{}k^l`$ in the potential $`V_{l^{},l}(k^{},k)`$.
The OPE potential in the singlet $`{}_{}{}^{1}S_{0}^{}`$ is regular at the origin and hence one can take $`\alpha _0`$ as an independent parameter or not (see Refs. Pavon Valderrama and Ruiz Arriola (2004a, b).) Actually, the smallness of the scattering length for the regular solution, suggests using the irregular solution. In the Weinbergโs counting of the potential, at NLO one has TPE contributions in the potential. At short distances they behave as an attractive $`1/r^5`$ potential (see Sect. VIII), and then $`\alpha _0`$ must be an independent parameter. At NNLO one has, again, a singular attractive $`1/r^6`$ potential (see Sect. IV), and thus an adjustable scattering length. This looks quite natural because increasing the order in the potential has a meaning and we can always compare the effect in the phase shifts of having a higher order potential with the same scattering length (See Sect. VIII). In this construction, if the next term in the expansion turned out to be more singular and repulsive the scattering length would be fully predicted from the potential.
The OPE potential in the triplet $`{}_{}{}^{3}S_{1}^{}^3D_1`$ coupled channel corresponds to case 2) and hence one has one free parameter in addition to the OPE potential parameters. One may choose this parameter to be the deuteron binding energy (or alternatively the triplet S-wave scattering length). Any other bound state or scattering observables are predicted. This case was treated in great detail in our previous work Pavon Valderrama and Ruiz Arriola (2005a) <sup>9</sup><sup>9</sup>9Relevant previous work on this channel was also presented in Ref. Beane et al. (2002) and Pavon Valderrama and Ruiz Arriola (2004b) where the orthogonality conditions where not considered. See the discussion at the end of Sect. V in Ref. Pavon Valderrama and Ruiz Arriola (2005a) . In the NLO TPE potential we have case 3) because both eigen potentials present a repulsive $`1/r^5`$ singularity (see Sect. VIII) and one would predict all observables from the potential parameters. Finally, in the NNLO TPE potential we have case 1) corresponding to an attractive-attractive (see Sect. V) and two additional parameters need to be specified for a state with a given energy. The orthogonality condition imposes a relation between two states of different energy, so that for all energies in the triplet channel one has three independent parameters. We will take these three parameters to be the deuteron binding energy, the asymptotic D/S ratio and the $`S`$wave scattering length. The trend one observes when going from LO to NNLO is quite natural; as usual in ChPT one has more parameters at any order of the approximation. The NLO approximation poses, however, a problem since one seems to have more predictive power than at LO (See Sect. VIII for more details on the consequences of using our renormalization ideas literally for the conventional NLO potential).
For the NNLO TPE triplet $`{}_{}{}^{3}S_{1}^{}^3D_1`$ channel we fix the deuteron binding energy, or equivalently $`\gamma `$, and the asymptotic D/S ratio $`\eta `$ by their experimental values. This fixes the short distance phases $`\phi _+(\gamma )`$ and $`\phi _{}(\gamma )`$. Next, if we use an $`\alpha `$ or $`\beta `$ (see below for a definition) zero energy scattering state we have in principle two short distance phases $`(\phi _{\alpha ,+}(0),\phi _{\alpha ,}(0))`$ and $`(\phi _{\beta ,+}(0)`$, $`\phi _{\beta ,}(0))`$ which can be related to the $`(\alpha _0,\alpha _{02})`$ and $`(\alpha _{02},\alpha _2)`$ scattering lengths respectively. Using the orthogonality constraints to the deuteron bound state one can then eliminate $`\alpha _{02}`$ and $`\alpha _2`$ and treat $`\alpha _0`$ as a free parameter. Thus, in the triplet $`{}_{}{}^{3}S_{1}^{}^3D_1`$ channel we can treat $`\gamma `$, $`\eta `$ and $`\alpha _0`$ as independent parameters. Once these parameters have been fixed we can actually predict the corresponding phase shifts since any positive energy state must be orthogonal both to the deuteron bound state and the zero energy scattering states. This result is a direct consequence of the singular Van der Waals attractive behavior of the TPE potential at the origin. It is remarkable that this same set of independent parameters was also adopted in Ref. de Swart et al. (1995) within the realistic potential model treatment.
Conflicts between naive dimensional power counting and renormalization have been reported recently already at the LO (OPE) level Nogga et al. (2005) where it is shown that even the $`{}_{}{}^{3}P_{0}^{}`$ partial wave depends strongly on the cut-off in momentum space (a gaussian regulator is used) if according to the standard Weinberg counting no counterterm is added. The requirement of renormalizability makes the promotion of one counter term unavoidable for channels which present an attractive singularity. This promotion is the minimal possible one compatible with finiteness, because in a coupled channel problem one could think in general of three counterterms. From this viewpoint the choice of just one counterterm in triplet channels is a bit mysterious. Our discussion in coordinate space agrees with these authors in the OPE potential, and actually allows to identify a priori the necessarily promotable counterterms as non trivial boundary conditions at the origin for singular attractive potentials. Moreover, we also see that the promotion of only one counterterm in the triplet channels with an attractive-repulsive singularity invoked in Ref. Nogga et al. (2005) is also maximal, since any additional counterterm would also produce divergent results (see Sect. VIII.3). Thus, we see that although power counting determines the long distance potential, the short distance singular character of the potential does not allow to fix the counterterms arbitrarily.
To conclude this discussion, let us mention that the short distance phases, whenever they become relevant play the role of some dimensionless constants which depend exclusively on the form of the potential, but not on the potential parameters Pavon Valderrama and Ruiz Arriola (2005a). For the same reason they can be taken to be zeroth order in the power counting used to generate the chiral potential in Eq. (1), although they are subjected in general to higher order corrections. In this sense, the form of the short distance interaction is dictated by the potential only, and cannot be considered as independent information (See also Ref. Griesshammer (2005) for a similar view on the three body problem in the absence of long distance potentials).
## IV The Singlet $`{}_{}{}^{1}S_{0}^{}`$-Channel
For the singlet $`{}_{}{}^{1}S_{0}^{}`$ channel one has to solve
$`u^{\prime \prime }(r)+U_{{}_{}{}^{1}S_{0}^{}}(r)u(r)`$ $`=`$ $`k^2u(r)`$ (33)
At short distances the NNLO NN chiral potential behaves as Kaiser et al. (1997); Friar (1999); Rentmeester et al. (1999)
$`U_{{}_{}{}^{1}S_{0}^{}}(r)`$ $``$ $`{\displaystyle \frac{3g^2}{128f^4\pi ^2r^6}}(4+15g^2+24\overline{c}_38\overline{c}_4)`$ (34)
$`=`$ $`{\displaystyle \frac{R^4}{r^6}}`$
which is a Van der Waals type interaction with typical length scale $`R=(MC_6)^{1/4}`$. Here, $`\overline{c}_i=Mc_i`$. The value of the coefficient is negative for the four parameter sets of Table 1, so the solution at short distances is of oscillatory type, Eq. (LABEL:eq:uA) with $`n=6`$, and
$`u(r)`$ $``$ $`A\left({\displaystyle \frac{r}{R}}\right)^{3/2}\mathrm{sin}\left[{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{R}{r}}\right)^2+\phi \right]`$ (35)
where there is a undetermined energy independent phase, $`\phi `$, and $`A`$ is a normalization constant. Note that the corresponding Van der Waals radius is quite sensitive to the choice of chiral parameters.
### IV.1 Low energy parameters
For the zero energy state we use the asymptotic normalization at large distances
$`u_0(r)`$ $``$ $`1{\displaystyle \frac{r}{\alpha _0}}.`$ (36)
Then, the effective range is given by
$`r_0`$ $`=`$ $`2{\displaystyle _0^{\mathrm{}}}๐r\left[\left(1{\displaystyle \frac{r}{\alpha _0}}\right)^2u_0(r)^2\right]`$ (37)
We can use the superposition principle for boundary conditions
$`u_0(r)`$ $`=`$ $`u_{0,c}(r){\displaystyle \frac{1}{\alpha _0}}u_{0,s}(r)`$ (38)
where $`u_{0,c}(r)1`$ and $`u_{0,s}(r)r`$ correspond to cases where the scattering length is either infinity or zero respectively. Using this decomposition one gets
$`r_0`$ $`=`$ $`A+{\displaystyle \frac{B}{\alpha _0}}+{\displaystyle \frac{C}{\alpha _0^2}},`$ (39)
where
$`A`$ $`=`$ $`2{\displaystyle _0^{\mathrm{}}}๐r(1u_{0,c}^2),`$ (40)
$`B`$ $`=`$ $`4{\displaystyle _0^{\mathrm{}}}๐r(ru_{0,c}u_{0,s}),`$ (41)
$`C`$ $`=`$ $`2{\displaystyle _0^{\mathrm{}}}๐r(r^2u_{0,s}^2),`$ (42)
depend on the potential parameters only. The interesting thing is that all explicit dependence on the scattering length $`\alpha _0`$ is displayed by Eq. (39 ). In a sense this is the non-perturbative universal form of a low energy theorem, which applies to any potential regular or singular at the origin which decays faster than a certain power of $`r`$ at large distances (for an analytical example with the pure Van der Waals potential $`U=R^4/r^6`$ see Sect. VII). Since the potential is known accurately at long distances we can visualize Eq. (39) as a long distance correlation between $`r_0`$ and $`\alpha _0`$. Naturally, if there is scale separation between the different contributions in the potential, Eq. (1), we expect the coefficients A,B and C to display a converging pattern. This is exactly what happens (see Eq. (LABEL:eq:r0\_todos) and Eq. (89) below) although not compatible with a naive and perturbative power counting (see Sect. IX).
In the $`{}_{}{}^{1}S_{0}^{}`$ the TPE potential becomes singular and attractive at short distances. Nevertheless, already at this point one can see the dramatic difference between attractive and repulsive singular potentials. In the attractive case the short distance phase allows to choose the scattering length independently on the potential, hence the coefficients $`A`$,$`B`$ and $`C`$ are uncorrelated with $`\alpha _0`$. For a singular repulsive potential, however, $`\alpha _0`$ as well as $`A`$,$`B`$ and $`C`$ are determined by the potential. If one assumes $`A`$,$`B`$ and $`C`$ to be independent on $`\alpha _0`$ in the repulsive case, this can only be possible due to an admixture of both the regular and irregular solutions, the latter will dominate at short distances and the effective range will diverge $`r_0\mathrm{}`$. This fact will become relevant in Sect. VIII.3.
Obviously, for the chiral TPE potential, Eq. (1), the coefficients have to be evaluated by numerical means and they are finite. We expect that these coefficients scale with the relevant scale of the potential. If long distances dominate $`A1/m`$, $`B1/m^2`$ and $`C1/m^3`$ but then $`r_01/m`$. On the contrary if short distances dominate $`AR`$, $`BR^2`$ and $`CR^3`$ and $`r_0R`$. The real situation is somewhat in between, but it is clear that $`A`$ is far more sensitive to short distances than $`C`$. Actually, for a large scattering length, as it is the case in the $`{}_{}{}^{1}S_{0}^{}`$ channel, the coefficient $`A`$ dominates. Note that unlike the standard approaches, where a short distance contribution to the effective range is allowed (in the form of a momentum dependent counterterm $`C_2(k^2+k^2)`$ ), we build $`r_0`$ solely from the potential and the scattering length $`\alpha _0`$. This is a direct consequence of the orthogonality relations, which preclude energy dependent boundary conditions for the local and energy independent chiral TPE potential.
In Fig. (3) we show the dependence of the effective range as a function of the short distance cut-off radius $`r_c`$, i.e. replacing the lower limit of integration in Eq. (37), for values between $`2`$ and $`0.1\mathrm{fm}`$ and taking the experimental value of the scattering length $`\alpha _0=23.74\mathrm{fm}`$. As we see, the short distance behaviour is well under control and nicely convergent towards the experimental value. This dependence also illustrates that an error estimate based on varying the cut-off between certain range is only a measure on the size of finite cut-off effects, rather than a measure on the error. The linear behaviour observed at small $`r_c`$ is a consequence of the dominance of the first term in Eq. (37) as compared to the second term where the wave function contribution vanishes as $`r_c^4`$. Let us remind that in the conventional treatments a counterterm $`C_2`$ is added to provide a short distance contribution to the effective range parameter and the result is fitted to experiment so that $`r_0`$ becomes an input of the calculation. Obviously one expects $`C_2`$ to depend on the regularization scale. As shown in Fig. (3) the size of the counterterm $`C_2`$ at $`r_c0`$ must be numerically small, since the TPE potential provides the bulk of the contribution. This agrees with the orthogonality constraint which requires $`C_20`$ when $`r_c0`$.
Numerically we find the following renormalized relations in the singlet channel for the OPE and NNLO TPE,
$`r_0`$ $`=`$ $`1.308062{\displaystyle \frac{4.547741}{\alpha _0}}+{\displaystyle \frac{5.192606}{\alpha _0^2}}(\mathrm{OPE}),`$
$`r_0`$ $`=`$ $`2.670963{\displaystyle \frac{5.755234}{\alpha _0}}+{\displaystyle \frac{6.031119}{\alpha _0^2}}(\mathrm{Set}\mathrm{I}),`$
$`r_0`$ $`=`$ $`2.715075{\displaystyle \frac{5.847358}{\alpha _0}}+{\displaystyle \frac{6.093430}{\alpha _0^2}}(\mathrm{Set}\mathrm{II}),`$
$`r_0`$ $`=`$ $`2.586862{\displaystyle \frac{5.584383}{\alpha _0}}+{\displaystyle \frac{5.916900}{\alpha _0^2}}(\mathrm{Set}\mathrm{III}),`$
$`r_0`$ $`=`$ $`2.616830{\displaystyle \frac{5.640921}{\alpha _0}}+{\displaystyle \frac{5.952694}{\alpha _0^2}}(\mathrm{Set}\mathrm{IV})`$
As we see, the coefficient dependent of $`\alpha _0^2`$ is not very sensitive to the choice of the coefficients $`c_1`$, $`c_3`$, $`c_4`$ and the OPE potential already provides the bulk of the contribution. On the other hand, the coefficient independent on $`\alpha _0`$ changes dramatically when going from OPE to TPE, suggesting that the effect is clearly non-perturbative. A direct inspection of the integrands for the A,B and C coefficients shows that $`A`$ picks its main contribution from the short distance region around 1 fm, whereas for $`B`$ and $`C`$ the most important contribution is located around 3 fm. One expects that different choices of coefficients $`c_3`$ and $`c_4`$ influence mostly the A coefficient. We confirm this expectation analytically by only keeping the Van der Waals contribution to the full potential in Sect. VII. We emphasize, again, that $`A`$, $`B`$ and $`C`$ are intrinsic information of the potential; these values for the effective range stem solely from the NNLO chiral potential and the scattering length $`\alpha _0`$, without any additional short distance contribution. The closeness of these numbers to the experimental value suggests that there is perhaps no need to make the boundary condition energy dependent if the cut-off is indeed removed, and that the missing $`0.1\mathrm{fm}`$ contribution can be clearly attributed to N<sup>3</sup>LO contributions in the potential.
The results are summarized in Table 2. For $`pn`$ we have the experimental values $`\alpha _0=23.74(2)`$ and $`r_0=2.77(5)`$. The previous formula, Eq. (39) yields $`r_0=2.92,2.97,2.83,2.87`$ for Sets I,II, III and IV respectively, which show a systematic discrepancy with the published values in several works (see References at the Table 2) and also a systematic trend to discrepancy with respect to the experimental value. Our renormalized values are always larger than the finite cut-off results. This seems natural since finite cut-off corrections diminish the integration region. Note also that the size of the discrepancy is larger than the experimental uncertainties and hence is statistically significant.
The value of the effective range was not given in the coordinate space calculation of Ref. Rentmeester et al. (1999) but the quality of the fit suggests that they get a value very close to the experimental one, $`r_0=2.75`$. The contribution to the effective range from the origin to $`0.1\mathrm{fm}`$ is about $`0.2`$. In Ref. Rentmeester et al. (1999) the cut-off is in coordinate space and an energy dependent boundary condition is considered. This means in practice cutting-off the lower integration in Eq. (37) at $`a=1.4\mathrm{fm}`$ and adding a short distance contribution $`r_S`$ as to reproduce the experimental value. This introduces a new potential independent parameter. As we have argued, in the limit $`a0`$, the short distance contribution of the effective range should go to zero, as implied by the orthogonality constraints. For finite $`a`$, the orthogonality constraint does not imply a vanishing short distance contribution to the effective range.
For Set IV one could reach the upper experimental value by flipping the sign of $`c_1`$ and keeping $`c_3`$ and $`c_4`$ unchanged. For $`c_1=2.43\mathrm{GeV}^1`$ one gets $`r_0=2.78\mathrm{fm}`$. The full experimental range would be covered by letting $`0.81\mathrm{GeV}^1<c_1<4.90\mathrm{GeV}^1`$. This is in total contradiction to the expectations of $`\pi N`$ scattering studies Buettiker and Meissner (2000). The insensitivity of our results with respect to the $`c_1`$ coefficient has to do with the fact that $`c_1`$ only enters in the potential at short distances at order $`1/r^4`$ which is sub-leading as compared to the leading Van der Waals singularity. This is another confirmation on the short distance dominance in the effective range parameter $`r_0`$.
Thus, according to our analysis, for the accepted values of chiral constants of Sets I,II, III and IV used in previous works, the difference in the value of the $`{}_{}{}^{1}S_{0}^{}`$ effective range could only be attributed to three pion exchange, relativistic effects and electromagnetic corrections. Another possibility, of course, is to refit the chiral constants to our renormalized, cut-off free results. This will be discussed in Sect. VI.
### IV.2 Phase Shift
For a finite energy scattering state we solve for the chiral TPE potential with the normalization
$`u_k(r){\displaystyle \frac{\mathrm{sin}(kr+\delta _0)}{\mathrm{sin}\delta _0}}`$ (44)
Again, if we use the superposition principle
$`u_k(r)=u_{k,c}(r)+k\mathrm{cot}\delta _0u_{k,s}(r)`$ (45)
with $`u_{k,c}\mathrm{cos}(kr)`$ and $`u_{k,s}\mathrm{sin}(kr)/k`$ and impose the orthogonality constraint with the zero energy state to get
$`{\displaystyle \frac{u_{k,c}^{}(a)+k\mathrm{cot}\delta _0u_{k,s}^{}(a)}{u_{k,c}(a)+k\mathrm{cot}\delta _0u_{k,s}(a)}}={\displaystyle \frac{\alpha _0u_{0,c}^{}(a)+u_{0,s}^{}(a)}{\alpha _0u_{0,c}(a)+u_{0,s}(a)}}`$ (46)
Note that the dependence of the phase-shift on the scattering length is explicit; $`\mathrm{cot}\delta _0`$ is a bilinear rational mapping of $`\alpha _0`$. Taking the limit $`a0`$ we get
$`k\mathrm{cot}\delta _0={\displaystyle \frac{\alpha _0๐(k)(k)}{\alpha _0๐(k)๐(k)}}`$ (47)
whereas the functions $`๐`$, $``$, $`๐`$ and $`๐`$ are even functions of $`k`$ which depend only on the potential and are given by
$`๐(k)`$ $`=`$ $`\underset{a0}{lim}\left(u_{0,c}(a)u_{k,c}^{}(a)u_{0,c}^{}(a)u_{k,c}(a)\right)`$
$`(k)`$ $`=`$ $`\underset{a0}{lim}\left(u_{k,c}(a)u_{0,s}^{}(a)u_{0,s}(a)u_{k,c}^{}(a)\right)`$
$`๐(k)`$ $`=`$ $`\underset{a0}{lim}\left(u_{0,c}^{}(a)u_{k,s}(a)u_{0,c}(a)u_{k,s}^{}(a)\right)`$
$`๐(k)`$ $`=`$ $`\underset{a0}{lim}\left(u_{0,s}(a)u_{k,s}^{}(a)u_{0,s}^{}(a)u_{k,s}(a)\right)`$
The obvious conditions $`๐(0)=๐(0)=0`$ and $`(0)=๐(0)=1`$ are satisfied. Expanding the expression for small $`k`$ one gets the well known effective range expansion
$`k\mathrm{cot}\delta ={\displaystyle \frac{1}{\alpha _0}}+{\displaystyle \frac{1}{2}}r_0k^2+v_2k^2+\mathrm{}`$ (49)
where $`v_k`$ is a polynomial in $`1/\alpha _0`$ of degree $`k+1`$.
The renormalized phase shift is presented in Fig. 4 for Set IV. As we see the trend in the effective range $`r_0`$ and the $`v_2`$ parameter is reflected in the behavior of the phase shift <sup>10</sup><sup>10</sup>10Let us mention that the momentum space calculation of Ref. Nogga et al. (2005) does not reproduce the physical and well measured scattering length. We have checked by an explicit solution of the Lipmann-Schwinger equation that the problem is more related to an insufficient number of Gauss points rather than to the value of the cut-off.
## V The Triplet $`{}_{}{}^{3}S_{1}^{}^3D_1`$ Channel
The coupled channel $`{}_{}{}^{3}S_{1}^{}^3D_1`$ set of equations read
$`u^{\prime \prime }(r)+U_{{}_{}{}^{3}S_{1}^{}}(r)u(r)+U_{E_1}(r)w(r)`$ $`=`$ $`k^2u(r),`$
$`w^{\prime \prime }(r)+U_{E_1}(r)u(r)+\left[U_{{}_{}{}^{3}D_{1}^{}}(r)+{\displaystyle \frac{6}{r^2}}\right]w(r)`$ $`=`$ $`k^2w(r),`$
At short distances the NN chiral NNLO potential behaves as Kaiser et al. (1997); Friar (1999); Rentmeester et al. (1999)
$`U_{{}_{}{}^{3}S_{1}^{}}(r)`$ $``$ $`{\displaystyle \frac{MC_{6,^3S_1}}{r^6}}`$
$`U_{E_1}(r)`$ $``$ $`{\displaystyle \frac{MC_{6,E_1}}{r^6}}`$
$`U_{{}_{}{}^{3}D_{1}^{}}(r)`$ $``$ $`{\displaystyle \frac{MC_{6,^3D_1}}{r^6}}`$
which is a coupled channels Van der Waals type interaction where the coefficients are given by
$`MC_{{}_{}{}^{3}S_{1}^{}}`$ $`=`$ $`{\displaystyle \frac{3g^2}{128f^4\pi ^2}}(43g^2+24\overline{c}_38\overline{c}_4)`$
$`MC_{E_1}`$ $`=`$ $`{\displaystyle \frac{3\sqrt{2}g^2}{128f^4\pi ^2}}(4+3g^216\overline{c}_4)`$
$`MC_{{}_{}{}^{3}D_{1}^{}}`$ $`=`$ $`{\displaystyle \frac{9g^2}{32f^4\pi ^2}}(1+2g^2+2\overline{c}_32\overline{c}_4)`$
If we diagonalize the corresponding matrix we get
$`\left(\begin{array}{cc}C_{6,^3S_1}& C_{6,E_1}\\ C_{6,E_1}& C_{6,^3D_1}\end{array}\right)`$ $`=`$ $`\left(\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta \\ \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right)\left(\begin{array}{cc}C_{6,+}& 0\\ 0& C_{6,}\end{array}\right)`$
$`\times `$ $`\left(\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta \\ \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right)`$
where $`C_{6,\pm }`$ are the corresponding eigenvalues and $`\theta `$ the mixing angle. They are listed in Table 1 for different parameters choices of the chiral couplings $`c_1`$, $`c_3`$ and $`c_4`$. We see that in all cases both eigenpotentials are attractive at short distances and hence the short distance behavior of the wave functions is of oscillatory type with $`n=6`$. Defining the Van der Waals scales
$`R_\pm =(MC_{6,\pm })^{1/4}`$ (54)
the short distance solutions read
$`\left(\begin{array}{c}u\\ w\end{array}\right)`$ $``$ $`\left(\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta \\ \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right)\left(\begin{array}{c}\left(\frac{r}{R_+}\right)^{\frac{3}{2}}\mathrm{sin}\left[\frac{1}{2}\left(\frac{R_+}{r}\right)^2+\phi _+\right]\\ \left(\frac{r}{R_{}}\right)^{\frac{3}{2}}\mathrm{sin}\left[\frac{1}{2}\left(\frac{R_{}}{r}\right)^2+\phi _{}\right]\end{array}\right)`$
Thus, we have two arbitrary short distance phases $`\phi _\pm `$ for a given fixed energy which cannot be deduced from the potential and hence have to be treated as independent parameters. We will fix them to some physical observables by integrating Eqs. (LABEL:eq:sch\_coupled) from infinity down to the origin.
### V.1 The deuteron
In the deuteron $`k^2=\gamma ^2`$ and we solve Eq. (LABEL:eq:sch\_coupled) together with the asymptotic condition at infinity
$`u(r)`$ $``$ $`A_Se^{\gamma r},`$
$`w(r)`$ $``$ $`A_De^{\gamma r}\left(1+{\displaystyle \frac{3}{\gamma r}}+{\displaystyle \frac{3}{(\gamma r)^2}}\right),`$ (55)
where $`\gamma =\sqrt{MB}`$ is the deuteron wave number, $`A_S`$ is the normalization factor and the asymptotic D/S ratio parameter is defined by $`\eta =A_D/A_S`$. In what follows we use $`\gamma `$ and $`\eta `$ as input parameters thus fixing the short distance phases $`\phi _\pm `$ automatically.
In this paper we compute the matter radius, which reads,
$`r_m^2={\displaystyle \frac{r^2}{4}}={\displaystyle \frac{1}{4}}{\displaystyle _0^{\mathrm{}}}r^2(u(r)^2+w(r)^2)๐r,`$ (56)
the quadrupole moment (without meson exchange currents)
$`Q_d={\displaystyle \frac{1}{20}}{\displaystyle _0^{\mathrm{}}}r^2w(r)(2\sqrt{2}u(r)w(r))๐r,`$ (57)
the deuteron inverse radius
$`r^1={\displaystyle _0^{\mathrm{}}}๐r{\displaystyle \frac{u(r)^2+w(r)^2}{r}},`$ (58)
which appears in low energy pion-deuteron scattering, and the $`D`$-state probability
$`P_D={\displaystyle _0^{\mathrm{}}}w(r)^2๐r.`$ (59)
Following Ref. Pavon Valderrama and Ruiz Arriola (2005a) we use the superposition principle of boundary conditions and write
$`u(r)`$ $`=`$ $`u_S(r)+\eta u_D(r),`$
$`w(r)`$ $`=`$ $`w_S(r)+\eta w_D(r),`$ (60)
where $`(u_S,w_S)`$ and $`(u_D,w_D)`$ correspond to the boundary conditions at infinity, Eq. (55), with $`A_S=1`$ and $`A_D=0`$ and with $`A_S=0`$ and $`A_D=1`$ respectively. Obviously, $`u_S,u_D,w_S`$ and $`w_D`$ depend on the potential and the deuteron binding energy only, so that the dependence on the asymptotic D/S ratio $`\eta `$ can de determined analytically. The value is taken as a free parameter. The resulting deuteron wave functions for Set IV are displayed in Fig. 5 and compared to the Nijmegen II results Stoks et al. (1993, 1994). One clearly sees the incommensurable ever increasing oscillations already below $`r=0.6\mathrm{fm}.`$
The short distance cut-off dependence of these deuteron properties using the experimental values for the deuteron binding energies and the asymptotic D/S ratio, $`\eta =0.0256`$, can be looked up in Fig. 6. As one sees the cut-off dependence is well under control, so the infinite cut-off limit can be extracted without difficulty.
Using the superposition principle of boundary conditions, Eq. (60 ) the asymptotic S-wave ratio depends quadratically on $`\eta `$ as follows
$`{\displaystyle \frac{1}{A_S^2}}`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}๐r(u_S^2+w_S^2)+2\eta {\displaystyle _0^{\mathrm{}}}๐r(u_Su_D+w_Sw_D)`$ (61)
$`+`$ $`\eta ^2{\displaystyle _0^{\mathrm{}}}๐r(u_D^2+w_D^2).`$
The coefficients of this second order polynomial depends on the potential and the deuteron binding energy. Similar relations hold for other observables. Evaluating the integrals numerically we get the following analytic correlations
Set I
$`1/A_S^2`$ $`=`$ $`3.78888214.675\eta +4489.43\eta ^2`$
$`r_m^2/A_S^2`$ $`=`$ $`5.4729754.1956\eta +1295.89\eta ^2`$
$`Q_d/A_S^2`$ $`=`$ $`0.342883+36.6449\eta 372.841\eta ^2`$
$`P_D/A_S^2`$ $`=`$ $`2.10904184.824\eta +4124.37\eta ^2`$
$`r^1/A_S^2`$ $`=`$ $`3.58173252.20\eta +5186.53\eta ^2`$ (62)
Set II
$`1/A_S^2`$ $`=`$ $`3.01271155.591\eta +3363.94\eta ^2`$
$`r_m^2/A_S^2`$ $`=`$ $`5.3473744.8896\eta +1122.59\eta ^2`$
$`Q_d/A_S^2`$ $`=`$ $`0.296852+33.2406\eta 309.624\eta ^2`$
$`P_D/A_S^2`$ $`=`$ $`1.44293132.314\eta +3098.89\eta ^2`$
$`r^1/A_S^2`$ $`=`$ $`2.52815171.36\eta +3635.18\eta ^2`$ (63)
Set III
$`1/A_S^2`$ $`=`$ $`4.65049283.545\eta +5902.53\eta ^2`$
$`r_m^2/A_S^2`$ $`=`$ $`5.5892963.1854\eta +1481.50\eta ^2`$
$`Q_d/A_S^2`$ $`=`$ $`0.377779+39.2691\eta 420.250\eta ^2`$
$`P_D/A_S^2`$ $`=`$ $`2.77521241.491\eta +5330.67\eta ^2`$
$`r^1/A_S^2`$ $`=`$ $`4.87639356.21\eta +7311.77\eta ^2`$ (64)
Set IV
$`1/A_S^2`$ $`=`$ $`3.40962190.713\eta +4198.86\eta ^2`$
$`r_m^2/A_S^2`$ $`=`$ $`5.4006649.2912\eta +1232.81\eta ^2`$
$`Q_d/A_S^2`$ $`=`$ $`0.306469+33.9354\eta 318.598\eta ^2`$
$`P_D/A_S^2`$ $`=`$ $`1.66525155.233\eta +3681.89\eta ^2`$
$`r^1/A_S^2`$ $`=`$ $`3.14658226.20\eta +4907.35\eta ^2`$ (65)
The numerical coefficients in these expressions depend on the deuteron binding energy and the TPE potential parameters, $`g`$, $`m`$, $`f`$, $`c_1`$, $`c_3`$ and $`c_4`$. The results for the deuteron properties are given in Table 3. The uncertainties are due to changing the input $`\gamma `$ and $`\eta `$ within their experimental uncertainties. We have checked that the short distance cut-offs $`a0.10.2\mathrm{fm}`$ generates much smaller uncertainties. The explicit dependence on $`\eta `$ is displayed in Fig. 7. Again, we find a discrepancy in the case of Set I with the values quoted in the finite cut-off calculation. Remarkably, our renormalized results in coordinate space agree most with the momentum space calculation of Ref. Entem and Machleidt (2002a) corresponding to Set IV. It is noticeable that this can be done without explicit knowledge of the counterterms used in that work in momentum space. This is precisely one of the points of renormalization; results can be reproduced by just providing physical input data, and no particular reference to the method of solution. Let us remind that the $`c_1`$, $`c_3`$ and $`c_4`$ were fixed from the perturbative study of NN peripheral waves where the cut-off sensitivity is rather small. Nevertheless, some significant discrepancies do also occur.
For the parameter Set IV Entem and Machleidt (2003b) obtained by a N<sup>3</sup>LO fit to NN scattering data, our NNLO calculation almost reproduces exactly the numbers provided in that work. Furthermore, they turn out to be compatible with the experimental numbers at the $`1\sigma `$ level within the uncertainty induced by the asymptotic $`D/S`$ ratio <sup>11</sup><sup>11</sup>11One may object that one should not use N<sup>3</sup>LO parameters to do a NNLO calculation, since they are obtained by fitting the same database. However, if there are finite cut-off effects the situation is not as clear. Finite cut-off effects are minimized in a N<sup>3</sup>LO calculation as shown in Ref. Epelbaum et al. (2005) where the induced uncertainties are drastically reduced when going from NNLO to N<sup>3</sup>LO. Note that in our calculation there are no sizeable cut-off induced uncertainties already at NNLO.
One immediate lesson we learn from inspection of Table 3 is that, regardless of the parameter set, only the experimental uncertainty in the asymptotic D/S ratio for the deuteron generates theoretical uncertainties about an order of magnitude larger then the experimental ones. On top of this, one has also to take into account other uncertainties, such as the one in $`g_{\pi NN}`$ and, of course, those induced by $`c_1`$, $`c_3`$ and $`c_4`$, which generally will generate larger uncertainties if all these parameters are regarded as independent (see Sect. VI below). In addition, there are systematic errors related to the accuracy of the expansion in the potential, Eq. (1). In common with non-perturbative finite cut-off calculations Rentmeester et al. (1999); Epelbaum et al. (2000); Epelbaum et al. (2004a, b); Entem and Machleidt (2003b); Rentmeester et al. (2003) they are difficult to estimate a priori given the non-perturbative nature of our calculation, but are bound to increase the error (see, however, our discussion in Sect. IX below on non-integer power counting). Given the insensitivity of our results with the short distance cut-off, the procedure used in Ref. Epelbaum et al. (2004a, b) of varying the cut-off becomes unsuitable in our case.
For the deuteron channel one may conclude that the predictive power of the chiral expansion has reached a limit at NNLO. So, at present, we do not expect to make theoretical predictions in the deuteron to be more accurate than experiment. The inclusion of N<sup>3</sup>LO and higher orders may provide better central values but is unlikely to improve the situation regarding error estimates since new unknown coefficients in the potential appear and the induced uncertainties will generally increase.
On the other hand, the slope for $`A_S`$ and $`r_m`$, Fig. 7, suggests that it would be better to take the asymptotic S-wave normalization or the matter radius as input, since generated errors may be comparable or even smaller. For instance, if the matter radius $`r_m`$ is taken as input we get instead $`\eta =0.0253(4)`$ a compatible value with similar errors. However, if we take $`A_S=0.8846(9)`$ as input for Set IV we get $`\eta =0.0255(1)`$ a compatible value with the experimental one but with much smaller errors. The reduction of errors is also confirmed in Sets I, II and III, although the central values are a bit off. This result opens up the possibility of making a benchmark determination of the asymptotic $`D/S`$ deuteron ratio from the chiral effective theory. Obviously, to do so, the chiral constants should be known with rather high accuracy, an illusory expectation at the present moment. In this regard it would perhaps be profitable to pin down the errors for the chiral constants from peripheral waves. This point will be analyzed elsewhere Pavon Valderrama and Ruiz Arriola (2006).
Both the loss of predictive power and the very rare possibility of making model independent theoretical predictions for purely hadronic processes using Chiral Perturbation Theory more accurate than experiment we seem to observe in low energy NN scattering is not new and has already been documented for low energy $`\pi \pi `$ scattering Nieves and Ruiz Arriola (2000); Colangelo et al. (2000, 2001) and provides a further motivation to use chiral effective approaches.
### V.2 Low energy parameters
The zero energy wave functions are taken asymptotically as <sup>12</sup><sup>12</sup>12We correct an error in Eq.(45) of our previous work Pavon Valderrama and Ruiz Arriola (2005a) where $`\alpha _2`$ appears. The corrected numerical value is $`\alpha _2=6.693\mathrm{fm}^5`$.
$`u_{0,\alpha }(r)`$ $``$ $`1{\displaystyle \frac{r}{\alpha _0}},`$
$`w_{0,\alpha }(r)`$ $``$ $`{\displaystyle \frac{3\alpha _{02}}{\alpha _0r^2}},`$
$`u_{0,\beta }(r)`$ $``$ $`{\displaystyle \frac{r}{\alpha _0}},`$
$`w_{0,\beta }(r)`$ $``$ $`\left({\displaystyle \frac{\alpha _2}{\alpha _{02}}}+{\displaystyle \frac{\alpha _{02}}{\alpha _0}}\right){\displaystyle \frac{3}{r^2}}{\displaystyle \frac{r^3}{15\alpha _{02}}}.`$ (66)
Using these zero energy solutions one can determine the effective range. The $`{}_{}{}^{3}S_{1}^{}`$ effective range parameter is given by
$`r_0`$ $`=`$ $`2{\displaystyle _0^{\mathrm{}}}\left[\left(1{\displaystyle \frac{r}{\alpha _0}}\right)^2u_\alpha (r)^2w_\alpha (r)^2\right]๐r.`$
(67)
In the zero energy case, the vanishing of the diverging exponentials at the origin imposes a condition on the $`\alpha `$ and $`\beta `$ states which generate a correlation between $`\alpha _0`$ , $`\alpha _{02}`$ and $`\alpha _2`$. Using the superposition principle of boundary conditions we may write the solutions in such a way that
$`u_{0,\alpha }(r)`$ $`=`$ $`u_1(r){\displaystyle \frac{1}{\alpha _0}}u_2(r)+{\displaystyle \frac{3\alpha _{02}}{\alpha _0}}u_3(r)`$
$`w_{0,\alpha }(r)`$ $`=`$ $`w_1(r){\displaystyle \frac{1}{\alpha _0}}w_2(r)+{\displaystyle \frac{3\alpha _{02}}{\alpha _0}}w_3(r)`$
$`u_{0,\beta }(r)`$ $`=`$ $`{\displaystyle \frac{1}{\alpha _0}}u_2(r)+\left({\displaystyle \frac{3\alpha _2}{\alpha _{02}}}+{\displaystyle \frac{3\alpha _{02}}{\alpha _0}}\right)u_3(r){\displaystyle \frac{1}{15\alpha _{02}}}u_4(r)`$
$`w_{0,\beta }(r)`$ $`=`$ $`{\displaystyle \frac{1}{\alpha _0}}w_2(r)+\left({\displaystyle \frac{3\alpha _2}{\alpha _{02}}}+{\displaystyle \frac{3\alpha _{02}}{\alpha _0}}\right)w_3(r){\displaystyle \frac{1}{15\alpha _{02}}}w_4(r)`$
where the functions $`u_{1,2,3,4}`$ and $`w_{1,2,3,4}`$ are independent on $`\alpha _0`$, $`\alpha _{02}`$ and $`\alpha _2`$ and fulfill suitable boundary conditions. The orthogonality constraints for the $`\alpha `$ and $`\beta `$ states read in this case
$`u_\gamma u_{0,\alpha }^{}u_\gamma ^{}u_{0,\alpha }+w_\gamma w_{0,\alpha }^{}w_\gamma ^{}u_{0,\alpha }|_{r=r_c}`$ $`=`$ $`0`$
$`u_\gamma u_{0,\beta }^{}u_\gamma ^{}u_{0,\beta }+w_\gamma w_{0,\beta }^{}w_\gamma ^{}u_{0,\beta }|_{r=r_c}`$ $`=`$ $`0`$
yielding two relations between $`\gamma `$, $`\alpha _{02}`$, $`\alpha _2`$, $`\eta `$ and $`\alpha _0`$, meaning that two of them are not independent. Using the superposition principle decomposition of the bound state, Eq. (60), and for the zero energy states, Eq. (LABEL:eq:sup\_zero), we make the orthogonality relation explicit in $`\alpha _0`$, $`\alpha _{02}`$, $`\alpha _2`$ and $`\eta `$. If we would use $`\alpha _0`$, $`\alpha _{02}`$, $`\alpha _2`$ as input parameters the orthogonality constraint is actually a non-linear eigenvalue problem for $`\gamma `$ and $`\eta `$. The values of $`\alpha _{02}`$ and $`\alpha _2`$ are not so well known although they have been determined in potential models in our previous work Pavon Valderrama and Ruiz Arriola (2004c). In contrast, $`\gamma `$, $`\eta `$ and $`\alpha _0`$ are well determined experimentally. Thus, in the deuteron scattering channel we will use $`\gamma `$, $`\eta `$ and $`\alpha _0`$ as independent input parameters and $`\alpha _{02}`$, $`\alpha _2`$ as predictions. This same set of independent parameters was also adopted in Ref. de Swart et al. (1995) within the high quality potential model treatment, although the role of the short distance Van der Waals singularity was not recognized. Fixing the experimental value of $`\gamma `$ we get the following relations for different parameter choices of $`c_1`$, $`c_3`$ and $`c_4`$, Set I
$`\alpha _{02}`$ $`=`$ $`{\displaystyle \frac{2.017630.456461\alpha _044.8947\eta +11.9351\alpha _0\eta }{0.314426+13.1555\eta }}`$
$`\alpha _2`$ $`=`$ $`{\displaystyle \frac{0.023522+1.04677\eta 11.6459\eta ^2+\alpha _0\left(0.0084230.537856\eta +9.39376\eta ^2\right)}{\alpha _0\left(0.023901+\eta \right)^2}}+{\displaystyle \frac{\alpha _{02}^2}{\alpha _0}}`$ (70)
Set II
$`\alpha _{02}`$ $`=`$ $`{\displaystyle \frac{1.717450.373228\alpha _033.4616\eta +8.76639\alpha _0\eta }{0.228075+9.865911\eta }}`$
$`\alpha _2`$ $`=`$ $`{\displaystyle \frac{0.030303+1.18083\eta 11.5032\eta ^2+\alpha _0\left(0.0095590.566611\eta +9.71850\eta ^2\right)}{\alpha _0\left(0.023118+\eta \right)^2}}+{\displaystyle \frac{\alpha _{02}^2}{\alpha _0}}`$ (71)
Set III
$`\alpha _{02}`$ $`=`$ $`{\displaystyle \frac{2.366590.550871\alpha _059.1666\eta +15.7488\alpha _0\eta }{0.414962+17.2806\eta }}`$
$`\alpha _2`$ $`=`$ $`{\displaystyle \frac{0.018755+0.937802\eta 11.7229\eta ^2+\alpha _0\left(0.0074370.505434\eta +9.09925\eta ^2\right)}{\alpha _0\left(0.024013+\eta \right)^2}}+{\displaystyle \frac{\alpha _{02}^2}{\alpha _0}}`$ (72)
Set IV
$`\alpha _{02}`$ $`=`$ $`{\displaystyle \frac{1.895260.418953\alpha _041.8369\eta +10.8526\alpha _0\eta }{0.279236+12.2978\eta }}`$
$`\alpha _2`$ $`=`$ $`{\displaystyle \frac{0.023751+1.04857\eta 11.5733\eta ^2+\alpha _0\left(0.0080500.518604\eta +9.31798\eta ^2\right)}{\alpha _0\left(0.023118+\eta \right)^2}}+{\displaystyle \frac{\alpha _{02}^2}{\alpha _0}}`$ (73)
The numerical coefficients appearing in these equations depend on the deuteron wave number $`\gamma `$ and the TPE parameters, $`g`$,$`f`$,$`m`$ and $`c_1`$, $`c_3`$ and $`c_4`$. The dependence on $`\eta `$ for fixed values of $`\alpha _0`$ within its experimental uncertainty is depicted in Fig. 8. We see that for fixed chiral couplings $`c_1`$, $`c_3`$ and $`c_4`$, the $`\eta `$ uncertainty dominates the errors. Numerical values can be seen at Table 3. Note the large discrepancy in the effective range $`r_0`$ for Sets I and II with the experimental number. Finite cut-off effects are observed in Set III although the $`\eta `$ induced uncertainty would make the value compatible with that estimate. Good agreement is observed again for Set IV, particularly in the $`E_1`$ and $`{}_{}{}^{3}D_{1}^{}`$ scattering lengths and the effective range $`r_0`$, but only the latter provides a clear TPE improvement as compared to OPE. The quantities $`\alpha _{02}`$ and $`\alpha _2`$ are compatible with typical expectations Pavon Valderrama and Ruiz Arriola (2004c) from the high quality potential models.
### V.3 Phase Shifts
For the $`\alpha `$ and $`\beta `$ positive energy scattering states we choose the asymptotic normalization
$`u_{k,\alpha }(r)`$ $``$ $`{\displaystyle \frac{\mathrm{cos}ฯต}{\mathrm{sin}\delta _1}}\left(\widehat{j}_0(kr)\mathrm{cos}\delta _1\widehat{y}_0(kr)\mathrm{sin}\delta _1\right),`$
$`w_{k,\alpha }(r)`$ $``$ $`{\displaystyle \frac{\mathrm{sin}ฯต}{\mathrm{sin}\delta _1}}\left(\widehat{j}_2(kr)\widehat{y}_2(kr)\mathrm{sin}\delta _1\right),`$
$`u_{k,\beta }(r)`$ $``$ $`{\displaystyle \frac{1}{\mathrm{sin}\delta _1}}\left(\widehat{j}_0(kr)\mathrm{cos}\delta _2\widehat{y}_0(kr)\mathrm{sin}\delta _2\right),`$
$`w_{k,\beta }(r)`$ $``$ $`{\displaystyle \frac{\mathrm{tan}ฯต}{\mathrm{sin}\delta _1}}\left(\widehat{j}_2(kr)\mathrm{cos}\delta _2\widehat{y}_2(kr)\mathrm{sin}\delta _2\right),`$
where $`\widehat{j}_l(x)=xj_l(x)`$ and $`\widehat{y}_l(x)=xy_l(x)`$ are the reduced spherical Bessel functions and $`\delta _1`$ and $`\delta _2`$ are the eigen-phases in the $`{}_{}{}^{3}S_{1}^{}`$ and $`{}_{}{}^{3}D_{1}^{}`$ channels, and $`ฯต`$ is the mixing angle $`E_1`$. The use of the superposition principle for boundary conditions as well as the orthogonality constraints,
$`u_\gamma u_{k,\alpha }^{}u_\gamma ^{}u_{k,\alpha }+w_\gamma w_{k,\alpha }^{}w_\gamma ^{}u_{k,\alpha }|_{r=r_c}`$ $`=`$ $`0`$
$`u_\gamma u_{k,\beta }^{}u_\gamma ^{}u_{k,\beta }+w_\gamma w_{k,\beta }^{}w_\gamma ^{}u_{k,\beta }|_{r=r_c}`$ $`=`$ $`0`$
analogous to Eq. (LABEL:eq:orth\_triplet), to the deuteron wave functions. If orthogonality would be applied to the zero energy state one obtains an explicit relation of $`\delta _1`$, $`\delta _2`$ and $`ฯต`$ with the scattering lengths $`\alpha _0`$, $`\alpha _2`$ and $`\alpha _{02}`$ as a direct generalization to the coupled channel case the one channel singlet case given by Eq. (47). The explicit expressions are rather cumbersome and will not be written down here explicitly. The results are depicted in Fig. 9 for Set IV. We observe a clear improvement in the threshold region, in consonance with the low energy parameters of Table 3 and a moderate improvement over the OPE results in the intermediate energy region. This suggests that finite cut-off effects may also be built in the phase shifts as well as the low energy parameters.
## VI Error analysis and Determination of chiral couplings from low energy NN data and the deuteron
### VI.1 Propagating experimental errors in $`c_1`$, $`c_3`$ and $`c_4`$.
The results in the previous sections clearly show that deuteron and low energy scattering properties in the $`{}_{}{}^{1}S_{0}^{}`$ and $`{}_{}{}^{3}S_{1}^{}^3D_1`$ channels are sensitive to finite cut-off effects and also to the values of the chiral constants after removal of the cut-off. We will assume that the values for $`c_1`$, $`c_3`$ and $`c_4`$ are free of uncertainties. Then, Set IV provides the best description of triplet data but produces a slightly off value for the effective range in the singlet channel at the $`2\sigma `$ confidence level. Note that in the singlet case the theoretical prediction for $`r_0`$ does not have a large source of error as in the triplet case where uncertainties in $`\eta `$ dominate the error. On the contrary, Set III provides a compatible value for the effective range in the singlet channel but incompatible values for the triplet channel in $`A_S`$ and $`r_0`$ at the $`34\sigma `$ confidence level on the experimental side. Thus, on this basis we may reject Set III and accept Set IV.
To improve on this analysis, let us try to include some errors on the chiral coefficients. The $`\pi N`$ analysis of Ref. Buettiker and Meissner (2000) (Set I) and the $`NN`$ fit of Ref.Rentmeester et al. (1999) (Set II) yield some errors. Ref. Entem and Machleidt (2003b) (Set IV) does not quote errors but we will take the educated guess of a $`5\%`$ error for $`c_3`$ and a $`30\%`$ error for $`c_4`$ Machleidt (2005). We can propagate then by a Monte-Carlo simulation implementing also the errors in $`g_{\pi NN}=13.1\pm 0.1`$, $`\alpha _{0,s}=23.77\pm 0.05`$, $`\alpha _{0,t}=5.419\pm 0.007`$, $`\eta =0.0256\pm 0.0004`$. We assume for simplicity that all these quantities are fully uncorrelated. This will in general enhance the errors, as compared to the case where correlations in $`c_1`$, $`c_3`$ and $`c_4`$ with $`\pi N`$ would be taken into account. Perhaps, the best thing would be to consider a simultaneous analysis of both $`NN`$ and $`\pi N`$ low energy data to build in correlations. Obviously, we do not expect good central values for the observables judging from Table 3. But there is still the possibility of large error bars.
The outcoming distributions in the low energy and deuteron parameters are somewhat asymmetric. Actually, for a given set of $`c_1`$, $`c_3`$ and $`c_4`$ distributions we observe the appearance of upper bounds in the $`{}_{}{}^{3}S_{1}^{}`$ effective range, namely $`r_{0,t}1.79,1,75,1.81\mathrm{fm}`$ for Sets I,II and IV respectively where the out-coming distributions become more dense. The results of the error propagation are summarized in Table 4. Thus, we see that the values of the chiral coefficients deduced from low energy $`\pi N`$ Buettiker and Meissner (2000) are globally inconsistent, at the $`1\sigma `$ level, with the low energy NN threshold parameters after uncertainties are taken into account. The same remark applies to Set II Rentmeester et al. (1999). Again, the loss of predictive power becomes manifest for all the sets although Set IV provides the best central values and the smallest errors. The situation for the quadrupole moment is noteworthy since the difference to the potential value is attributed to Meson Exchange Currents (MEC) and relativistic effects, which provide a correction of about $`0.01\mathrm{fm}^2`$ (see Ref. in Ref. Entem and Machleidt (2002a) and also Ref. Phillips (2003)). As we see, this is about the size of the error deduced from our analysis. In the case of the deuteron matter radius the situation is even worse since MECโs contributions are much smaller $`0.003\mathrm{fm}`$ Phillips (2003) while our predicted errors are larger. It would be extremely interesting to reanalyze the problem with the present deuteron wave functions Nogga et al. (2006).
### VI.2 Determination of $`c_1`$, $`c_3`$ and $`c_4`$.
Another possibility is to attempt a direct fit to the data. The standard approach is to fit the partial waves to a NN database Stoks et al. (1993, 1994). The problem with such an approach is that, unfortunately, there is no error assignment on the phase shifts and hence a reliable assessment of errors cannot be made. Actually, besides the work of Ref. Rentmeester et al. (1999, 2003) where a full partial wave analysis was undertaken, other works Entem and Machleidt (2003b); Epelbaum et al. (2004b); Epelbaum et al. (2005) assume fixed values for $`c_1`$ , $`c_3`$ and $`c_4`$ without attempting any error analysis based on input uncertainties. But even if data for the phase shifts with errors were known one expects the quality of the fit to worsen as the energy is increased as we think that the chiral approach to NN interaction should work best at low energies. If the data were known with uniform uncertainty one would fit until $`\chi ^2/DOF`$ exceeds one providing an energy window. In such a fit all points are equally weighted while we know that the description at low energies, where the theory works best, will be compromised by the highest possible energy within such an energy window.
Along the previous line of reasoning we propose, instead, to fit directly the low energy threshold parameters which central values and errors are well known and widely accepted. The basic ingredient is to use purely hadronic information in the process to avoid any contamination due to electromagnetic effects. Specifically, we make a Monte-Carlo sampling of the input data assuming that as primary data they are gaussian distributed and uncorrelated. For any of the samples we make a $`\chi ^2`$ fit to the values $`r_{0,s}`$, $`r_{0,t}`$ and $`A_S`$, i.e. we minimize
$`\chi ^2=\left({\displaystyle \frac{r_{0,s}r_{0,s}^{\mathrm{exp}}}{\mathrm{\Delta }r_{0,s}}}\right)^2+\left({\displaystyle \frac{r_{0,t}r_{0,t}^{\mathrm{exp}}}{\mathrm{\Delta }r_{0,t}}}\right)^2+\left({\displaystyle \frac{A_SA_S^{\mathrm{exp}}}{\mathrm{\Delta }A_S}}\right)^2`$
and determine then the optimal values of $`c_1`$ , $`c_3`$ and $`c_4`$. We only accept values where $`\chi ^2<1`$, and the resulting distribution of chiral constants $`c_1`$, $`c_3`$ and $`c_4`$ is given in Figs. 10. As we see, there is a very strong, almost linear, correlation between $`c_3`$ and $`c_4`$. This can be easily understood in terms of the short distance dominance of the singlet effective range, since for the pure Van der Waals contribution, and in the limit of large scattering length $`r_{0,s}R=(MC_6)^{1/4}`$, with $`C_6`$ given in Eq. (34). Deviations from linearity are induced from the larger relative error of $`r_{0,s}`$ ($`1\%`$) as compared to $`r_{0,t}`$ and $`A_S`$ ($`0.1\%`$). This is different from the large scale partial phase-shift analysis obtained in Ref. Rentmeester et al. (2003) where a very small correlation between $`c_3`$ and $`c_4`$ of about 0.2 was found. We have checked that cutting-off data with decreasing values of $`\chi ^2`$, excludes the points where the distribution is sparse, so that the dense part indeed reflects the uncertainties in the input data. The fact that the three coefficients seem to be on a line is just a consequence of solving by minimization a system of three equations and three unknowns.
We use $`g_{\pi NN}=13.083`$, $`\alpha _{0,s}=23.77\pm 0.05`$, $`\alpha _{0,t}=5.419\pm 0.007`$, $`\eta =0.0256\pm 0.0004`$ and fit $`c_1`$ , $`c_3`$ and $`c_4`$ to the values $`r_{0,s}=2.77\pm 0.05`$, $`r_{0,t}=1.753\pm 0.08`$ and $`A_S=0.8846\pm 0.0009`$. Our final result for a sample with 125 points with $`\chi ^2<1`$ is
$`c_1`$ $`=`$ $`1.13_{0.04}^{+0.02}(\mathrm{stat})\mathrm{GeV}^1,`$
$`c_3`$ $`=`$ $`2.60_{0.23}^{+0.18}(\mathrm{stat})\mathrm{GeV}^1,`$
$`c_4`$ $`=`$ $`+3.40_{0.40}^{+0.25}(\mathrm{stat})\mathrm{GeV}^1.`$ (78)
The central value is the mean and the errors have been obtained by the standard method of excluding the $`16\%`$ left and right extreme values of the variables, so as to have $`68\%`$ confidence level between the upper and lower values. Cutting-off data with $`\chi ^2<0.5`$ does not change significantly the result.
At the $`2\sigma `$ level, our values for $`c_1`$, $`c_3`$ and $`c_4`$ are compatible with the analysis of low energy $`\pi N`$ scattering of Ref. Buettiker and Meissner (2000), $`c_1=0.81\pm 0.15`$, $`c_3=4.69\pm 1.34`$ and $`c_4=3.40\pm 0.04`$, but incompatible with the NN full partial wave analyses Rentmeester et al. (1999, 2003) where an energy dependent boundary condition at $`a=1.4\mathrm{fm}`$ was used. It is difficult to say whether other determinations for the chiral couplings based on NN scattering are incompatible with ours, since no error estimates have been provided.
### VI.3 Estimate of the systematic errors
As we have mentioned, any approach based on power counting of the potential cannot make an a priori estimate of the accuracy of the calculation. Nevertheless, we can have an idea by simply varying the input parameters.
At LO we may use either $`g_A=1.26`$ as input or $`g_{\pi NN}=13.1`$, since the difference is the Goldberger-Treiman discrepancy, which should be a higher order correction. The effect can be seen by comparing OPE with OPE in Table 3. When compared to the TPE result, for e.g. Set IV, the error is underestimated this way. At NNLO we use the same procedure in the TPE piece. Again, the difference should be higher orders. Numerically this is equivalent to include $`g_{\pi NN}=13.1\pm 0.1`$ and consider this systematic error as an statistical one, which has already been taken into account.
We can estimate the systematic error in the chiral constants by varying the input used to determine the cโs. To correlate the singlet and triplet channels we must keep $`r_{0,s}`$ and $`\alpha _{0,s}`$. So we can interchange the inputs $`A_S`$, $`r_{0,t}`$ with the outputs $`r_m`$, $`Q_d`$, $`\alpha _{02}`$ and $`\alpha _2`$. This yields a total of 15 possible combinations. Another question concerns the assessment an error to the fitted variables whenever there is no direct experimental quantity, since this choice weights the determination of the cโs. This is the case for $`\alpha _{02}`$ and $`\alpha _2`$, where we make the educated guess of taking the difference between the Reid93 and NijmII values as determined in Ref. Pavon Valderrama and Ruiz Arriola (2004c) as an estimate of the error. The situation with $`Q_d`$ is a bit special and we exclude it from the analysis <sup>13</sup><sup>13</sup>13The discrepancy of potential models to the experimental value $`0.01\mathrm{fm}^2`$, attributed to MECโs and relativistic effects Phillips (2003), is about two orders of magnitude larger than the error in the experimental number $`0.0003\mathrm{fm}^2`$ and the discrepancy between potential models $`0.0004\mathrm{fm}^2`$. It is not clear whether the discrepancy can be pinned down with similar errors Nogga et al. (2006).. The results are listed in Table 5. We see that this estimate on the systematic error provides a larger fluctuation than direct propagation of the input errors for $`c_1`$. Symmetrizing the errors we get
$`c_1`$ $`=`$ $`1.2\pm 0.2(\mathrm{syst})\mathrm{GeV}^1,`$
$`c_3`$ $`=`$ $`2.6\pm 0.1(\mathrm{syst})\mathrm{GeV}^1,`$
$`c_4`$ $`=`$ $`+3.3\pm 0.1(\mathrm{syst})\mathrm{GeV}^1.`$ (79)
If we attempt a fit to all observables assigning $`\mathrm{\Delta }Q_d=0.01\mathrm{fm}^2`$ and $`\mathrm{\Delta }\alpha _{02}=0.4\mathrm{fm}`$ we get $`c_1=0.9`$, $`c_3=2.71`$ and $`c_4=3.85`$ with a large $`\chi ^2/DOF=3`$ basically due to the small errors. Obviously a more realistic estimate of the errors would be desirable.
## VII The role of chiral Van der Waals forces
As we have pointed out, our approach is not the conventional one of adding short distance counterterms following a given a priori power counting regardless on the approximation where the long distance potential has been constructed. Instead, the potential power counting dictates the form of the short distance physics by demanding a finite limit when the regulator is removed. In order to stress the differences with previous approaches it is interesting to see how much of the phase shifts is determined from the short distance chiral potential without adding a short range contribution to the effective range. In the standard approach this can be achieved by adding a counterterm $`C_2`$ in the $`S`$wave channels. In Ref. Pavon Valderrama and Ruiz Arriola (2005a) we showed that both perturbatively and non-perturbatively the orthogonality constraints for the OPE potential imply $`C_2=0`$. Here we will see that the bulk of the $`S`$wave interaction can be explained mainly in terms of the chiral Van der Waals force when renormalization is carried out, without any additional short distance contribution or counterterm.
For a pure Van der Waals potential of the form
$`U={\displaystyle \frac{R^4}{r^6}},`$ (80)
the zero energy wave function can be analytically computed Frank et al. (1971) in terms of Bessel functions $`J_\nu (x)`$. Normalizing to the asymptotic form $`u_0(r)1r/\alpha _0`$ we get
$`u_0(r)`$ $`=`$ $`\mathrm{\Gamma }\left({\displaystyle \frac{5}{4}}\right)\sqrt{{\displaystyle \frac{2r}{R}}}J_{\frac{1}{4}}\left({\displaystyle \frac{R^2}{2r^2}}\right)`$ (81)
$``$ $`\mathrm{\Gamma }\left({\displaystyle \frac{3}{4}}\right)\sqrt{{\displaystyle \frac{rR}{2}}}J_{\frac{1}{4}}\left({\displaystyle \frac{R^2}{2r^2}}\right){\displaystyle \frac{1}{\alpha _0}}.`$
The effective range can also be computed analytically Gao (1998); Flambaum et al. (1999) from Eq. (37) yielding
$`r_0`$ $`=`$ $`{\displaystyle \frac{4R^2}{3\alpha _0}}+{\displaystyle \frac{4R^3\mathrm{\Gamma }(\frac{3}{4})^2}{3\alpha _0^2\pi }}+{\displaystyle \frac{16R\mathrm{\Gamma }(\frac{5}{4})^2}{3\pi }},`$ (82)
$`=`$ $`1.39473R{\displaystyle \frac{1.33333R^2}{\alpha _0}}+{\displaystyle \frac{0.637318R^3}{\alpha _0^2}},`$
in agreement with the general low energy theorem of Eq. (39). Taking the values of Table 1 for $`R=(MC_6)^{1/4}`$ one gets in the singlet $`{}_{}{}^{1}S_{0}^{}`$ channel
$`r_{0,s}`$ $`=`$ $`2.39811{\displaystyle \frac{3.9418}{\alpha _{0,s}}}+{\displaystyle \frac{3.23959}{\alpha _{0,s}^2}}(\mathrm{Set}\mathrm{I}),`$
$`r_{0,s}`$ $`=`$ $`2.49192{\displaystyle \frac{4.25624}{\alpha _{0,s}}}+{\displaystyle \frac{3.63486}{\alpha _{0,s}^2}}(\mathrm{Set}\mathrm{II}),`$
$`r_{0,s}`$ $`=`$ $`2.2227{\displaystyle \frac{3.8625}{\alpha _{0,s}}}+{\displaystyle \frac{2.57944}{\alpha _{0,s}^2}}(\mathrm{Set}\mathrm{III}),`$
$`r_{0,s}`$ $`=`$ $`2.29099{\displaystyle \frac{3.59753}{\alpha _{0,s}}}+{\displaystyle \frac{2.82459}{\alpha _{0,s}^2}}(\mathrm{Set}\mathrm{IV}).`$
The numerical agreement at the ten percent level of the $`\alpha _{0,s}`$ independent term with the full chiral TPE result, Eq. (LABEL:eq:r0\_todos)), is striking <sup>14</sup><sup>14</sup>14The formula (82) can also be used as a numerical test of the integration method and of the numerical solution of the differential equations. This is a non-trivial condition due to the rapid oscillations of the wave function at the origin. We have checked that it is accurately reproduced.. On the other hand, first order perturbation theory in the OPE potential yields (see Sect. A of Ref. Pavon Valderrama and Ruiz Arriola (2005a)) in the form of Eq. (39) the result
$`r_{0,s}`$ $`=`$ $`{\displaystyle \frac{g_{\pi NN}^2}{8M\pi }}\left(1{\displaystyle \frac{8}{3\alpha _{0,s}m}}+{\displaystyle \frac{2}{\alpha _{0,s}^2m^2}}\right),`$ (84)
$`=`$ $`1.4369{\displaystyle \frac{5.4789}{\alpha _{0,s}}}+{\displaystyle \frac{5.8758}{\alpha _{0,s}^2}}.`$
Note that the coefficient in $`1/\alpha _{0,s}^2`$ is slightly better described by the OPE perturbative value than the full OPE result (see Eq. (LABEL:eq:r0\_todos)), a not unreasonable result since this coefficient is sensitive to the longest range part of the interaction. Likewise, the bulk of the $`\alpha _0`$-independent coefficient is given just by the most singular contribution to the full chiral potential. As we see, for large scattering lengths the effective range scales with the Van der Waals singlet radius $`R_s=(MC_{6,^1S_0})^{1/4}`$ and not with the pion Compton wavelength $`1/m`$, confirming the dominance of the short distances singularity in the singlet channel.
For the triplet channel, the equation cannot be solved analytically, and the effective range has a correction due to the D-wave (see Eq. (67) ). Moreover, the scattering length is a factor five times smaller than in the singlet case, so that we do not expect in principle such a dramatic agreement. If we neglect the mixing with the D-wave and take the $`R_t=(MC_{6,^3S_1})^{1/4}`$ of Eq. (LABEL:eq:vdw\_triplet) we get
$`r_{0,t}`$ $`=`$ $`2.50174{\displaystyle \frac{4.28983}{\alpha _{0,t}}}+{\displaystyle \frac{3.67797}{\alpha _{0,t}^2}}(\mathrm{Set}\mathrm{I}),`$
$`r_{0,t}`$ $`=`$ $`2.58537{\displaystyle \frac{4.58143}{\alpha _{0,t}}}+{\displaystyle \frac{4.05928}{\alpha _{0,t}^2}}(\mathrm{Set}\mathrm{II}),`$
$`r_{0,t}`$ $`=`$ $`2.35089{\displaystyle \frac{3.78809}{\alpha _{0,t}}}+{\displaystyle \frac{3.05196}{\alpha _{0,t}^2}}(\mathrm{Set}\mathrm{III}),`$
$`r_{0,t}`$ $`=`$ $`2.40877{\displaystyle \frac{3.97691}{\alpha _{0,t}}}+{\displaystyle \frac{3.28297}{\alpha _{0,t}^2}}(\mathrm{Set}\mathrm{IV}),`$
which, using the triplet scattering length value, $`\alpha _{0,t}=5.42`$ yields $`r_{0,t}=1.83,1.87,1.75,1.78`$ respectively, in remarkable agreement with the experimental value. An estimate of the mixing effect can be made by using the largest van der Waals eigen radius $`R_+=(MC_{6,+})^{\frac{1}{4}}`$ obtained by diagonalizing the interaction at short distances. From Table 1, Eq. (82) and the experimental value of the scattering length we get $`r_0=2.00,2.07,1.95,2.05\mathrm{fm}`$ for Sets I,II,III and IV respectively, accounting for about $`85\%`$ percent of the full value. Instead, perturbation theory for OPE, Eq. (84) yields $`r_0=0.62\mathrm{fm}`$, and full OPE $`r_0=1.64`$. Actually, using the relation
$`MC_{6,^1S_0}MC_{6,^3S_1}`$ $`=`$ $`R_s^4R_t^4`$ (86)
$`=`$ $`{\displaystyle \frac{3g^2}{64\pi ^2f^4}}(49g^2)`$
we get an explicit correlation between $`\alpha _{0,s}`$, $`r_{0,s}`$, $`\alpha _{0,t}`$ and $`r_{0,t}`$ regardless on the numerical values of the chiral constants $`c_3`$ and $`c_4`$. In the range of physical parameters this looks like a linear correlation (see Fig. 12) between the singlet and triplet effective ranges. For $`r_{0,t}=1.75`$ one gets $`r_{0,s}=2.34`$.
To check further the dominance of chiral Van der Waals interactions, we plot in Fig. 11 the phase shifts for a variety of situations including the pure Van der Waals contributions, as well as the contribution of the NNLO only, which reduces to the previous case at short distances but decays exponentially as $`e^{2mr}`$ at long distances. The plots confirm, again, our estimations based on the pure Van der Waals potential of the effective range for the $`s`$waves, and this is the reason why the triplet $`s`$wave is better reproduced than the singlet case for Set IV. Obviously, by adjusting the effective range changing the chiral parameters $`c_3`$ and $`c_4`$ we could obtain a much better description of the data.
The results of this study show that the singularity of the chiral Van der Waals force is not a feature that should be avoided, but instead provides a very simple way to describe the scattering data for the $`s`$waves.
Finally, it is interesting to note, that central waves based on taking the chiral limit of the potential are less accurately described than the phase-shifts obtained from the pure Van der Waals contribution. In this limit the singlet $`{}_{}{}^{1}S_{0}^{}`$ channel contains in addition to the Van der Waals term a $`1/r^5`$ contribution stemming from the NLO TPE contribution. The triplet $`{}_{}{}^{3}S_{1}^{}^3D_1`$ TPE contribution has a similar structure in addition to the OPE tensor $`1/r^3`$ singular short distance contribution.
## VIII The TPE Potential at NLO: A missing link ?
In the previous sections we analyzed the renormalization of the NNLO potential. In this Section we analyze the NLO in the singlet $`{}_{}{}^{1}S_{0}^{}`$ and triplet $`{}_{}{}^{3}S_{1}^{}^3D_1`$ channels and the problem that arises in the latter. We argue that similar trends are observed in finite cut-off calculations. We also suggest several scenarios on how the problem may be overcome.
### VIII.1 Convergence in the singlet $`{}_{}{}^{1}S_{0}^{}`$ channel
In the singlet $`{}_{}{}^{1}S_{0}^{}`$ channel the potential at short distances behaves as Kaiser et al. (1997); Friar (1999); Rentmeester et al. (1999)
$`U_{{}_{}{}^{1}S_{0}^{}}{\displaystyle \frac{MC_{5,^1S_0}}{r^5}},`$ (87)
where
$`MC_{5,^1S_0}`$ $`=`$ $`{\displaystyle \frac{M(1+10g^259g^4)}{256\pi ^3f^4}}`$ (88)
The singlet coefficient is negative and, according to the discussion in Sect. III, one has an undetermined short distance phase which can be fixed by using the scattering length as input. The effective range in the singlet channel is given by
$`r_0=2.122{\displaystyle \frac{4.889}{\alpha _0}}+{\displaystyle \frac{5.499}{\alpha _0^2}}(\mathrm{NLO}),`$ (89)
which compared with the LO and NNLO results, Eq. (LABEL:eq:r0\_todos), shows a convergence rate. To show that this trend to convergence is not fortuitous we display in Table 6 the threshold parameters of the effective range expansion $`k\mathrm{cot}\delta =1/\alpha _0+r_0k^2/2+v_2k^4+v_3k^6+v_4k^8`$ depending on the terms kept in the expansion of the potential given by Eq. (1). As we see there is a clear trend to convergence, although the higher order threshold parameters display a slower convergence rate since they are increasingly sensitive to the shorter range regions. This trend is confirmed in Fig. 13 for the phase shift. Obviously, there is scale separation in the singlet potential, and higher order potentials although more singular at the origin yield contributions in the right direction.
### VIII.2 The problem in the triplet $`{}_{}{}^{3}S_{1}^{}^3D_1`$ channel
The triplet $`{}_{}{}^{3}S_{1}^{}^3D_1`$ potential at short distances has the behaviour Kaiser et al. (1997); Friar (1999); Rentmeester et al. (1999)
$`U_{{}_{}{}^{3}S_{1}^{}}(r)`$ $``$ $`{\displaystyle \frac{MC_{5,^3S_1}}{r^5}},`$
$`U_{E_1}(r)`$ $``$ $`{\displaystyle \frac{MC_{5,E_1}}{r^5}},`$
$`U_{{}_{}{}^{3}D_{1}^{}}(r)`$ $``$ $`{\displaystyle \frac{MC_{5,^3D_1}}{r^5}},`$ (90)
where
$`MC_{5,^3S_1}`$ $`=`$ $`{\displaystyle \frac{3M(110g^2+27g^4)}{256\pi ^3f^4}},`$
$`MC_{5,E_1}`$ $`=`$ $`{\displaystyle \frac{15Mg^4}{64\sqrt{2}\pi ^3f^4}}`$
$`MC_{5,^3D_1}`$ $`=`$ $`{\displaystyle \frac{3M(110g^2+37g^4)}{256\pi ^3f^4}},`$ (91)
On the other hand, the diagonalized triplet coefficients are
$`MC_{5,+}`$ $`=`$ $`{\displaystyle \frac{3M(110g^2+17g^4)}{256\pi ^3f^4}},`$
$`MC_{5,}`$ $`=`$ $`{\displaystyle \frac{3M(110g^2+47g^4)}{256\pi ^3f^4}},`$ (92)
and the mixing angle is given by $`\mathrm{tan}\theta =\sqrt{2}`$, differing by $`\pi `$ as compared to the OPE case Pavon Valderrama and Ruiz Arriola (2005a). For $`0.5356<g<0.8217`$ one would have an attractive-repulsive situation ( see Sect. III), as in the OPE case Pavon Valderrama and Ruiz Arriola (2005a) and in such a case one could take either the deuteron binding energy or the $`{}_{}{}^{3}S_{1}^{}`$ scattering length. However, for the physical value $`g=1.26`$ one has two short distance repulsive eigenchannels, and hence one must take the exponentially decaying regular solutions at the origin. Let us remind the according to Sect. III finite renormalized results can only be obtained by precisely choosing the regular solution at the origin. In this case there are no short distance phases, and the scattering lengths, as well as the phase shifts are completely determined from the potential. The (finite) renormalized results are depicted in Fig. 13. As we see, the singlet $`{}_{}{}^{1}S_{0}^{}`$ phase-shift shows a very reasonable trend, since NLO and improves on the LO, and it is improved by the NNLO potential. We remind that in the three cases the scattering length is exactly the same. However, not completely unexpectedly, the triplet channel results worsen the LO ones. In the next subsection we show that if, demanding the standard Weinberg counting requires the irregular solution at the origin, hence yielding to divergent renormalized results.
### VIII.3 Finite cut-offs and the Weinberg counting
The special status of the NLO calculation as compared to the LO and NNLO ones has been recognized in previous studies in momentum space Epelbaum et al. (2000) where regularization was implemented by using a sharp cut-off $`\mathrm{\Lambda }`$. As noted by these authors, the allowed cut-off variations at NLO are smaller ($`380600`$ MeV ) than at LO ($`700800`$ MeV) or NNLO ($`8001000`$) but the reasons have not been made clear. Let us focus on the triplet $`{}_{}{}^{3}S_{1}^{}^3D_1`$ channel. Within our coordinate space renormalization scheme this trend can be easily understood. At LO one fixes only one parameter, say $`\alpha _0`$, and because one has attractive and repulsive potentials at short distances the system will naturally be driven into the exponential regular solution at the origin. Obviously, if one would fix some other parameter independently, say $`r_0`$ (or equivalently using a counterterm $`C_2`$, and not the one predicted by the regular solutions, one would be driven instead to the irregular solution, not allowing to remove the cut-off in practice. In such a situation one would be forced to keep the cut-off finite at the scale where the repulsive core sets in. However, at LO the Weinberg power counting does not allow to fix this additional parameter and one can comfortably reach higher cut-off values. On the contrary, at NLO one has two repulsive eigenpotentials and one cannot fix any low energy parameter arbitrarily. Otherwise one would be attracted to the irregular solutions at short distances. On the other hand, they are attractive at long distances, so that one would expect a stability region where the potential becomes flat before turning into a repulsive core in both eigenchannels. This is exactly what one observes in the NLO calculation of Ref. Epelbaum et al. (2000). The occurrence of such a plateau is to some extent fortuitous since it is associated to the critical points of the potential, and not with some a priori estimate on the validity range of the NLO potential. Finally, in the NNLO calculation because both eigenpotentials have an attractive character one can again increase the cut-off since there are no irregular solutions in the problem where one can be attracted to. It is very rewarding that our coordinate space analysis of short distance singularities anticipates when these features can be expected. On the other hand this does not imply that finite cut-off calculations are necessarily wrong, simply that the observed features when the cut-off approaches the limit can be understood.
The previous discussion can be illustrated in our approach by looking at the short distance cut-off dependence of the effective range, $`r_0`$ (in fm) and the deuteron wave function renormalization, $`A_S`$ (in $`\mathrm{fm}^{1/2}`$), in the triplet $`{}_{}{}^{3}S_{1}^{}^3D_1`$ channel at LO, NLO and NNLO in the standard Weinberg counting as presented in Fig. 14. As described in Sect. III.3 any counterterms can be mapped into a given renormalization condition. Once these conditions are fixed we can ask whether other properties are finite or not. Thus, at LO we fix the deuteron binding energy (one counterterm), at NLO and NNLO we fix the deuteron biding energy, the asymptotic $`D/S`$-mixing, $`\eta `$, and the scattering length $`\alpha _0`$ (three counterterms). In all cases it is clear that by lowering the cut-off at LO and NNLO of the approximation one nicely approaches the experimental values. This raises immediately the question whether there is a given value of the cut-off where NLO improves over LO. As we see such a region does not exist. In addition, although there is a nice and clear trend in both LO and NNLO for distances below $`0.5\mathrm{fm}`$ for $`A_S`$ down to the origin, this is not so at NLO. So, in this case it is not true that low energy properties are independent on short distance details, in contrast to the standard EFT wisdom. Moreover, Fig. 14 shows explicitly the conflict between the Weinberg counting and the remotion of the cut-off at NLO because $`r_0`$ and $`A_S`$ diverge due to the onset of the irregular solution, as anticipated in our study of short distance solutions (see Sect. III). We have checked that this is a general feature on both deuteron and scattering properties. On the contrary, LO and NNLO have a rather smooth limit because in these two cases Weinberg power counting on the short distance counterterms turns out to be compatible with the choice of the regular solution at the origin. Thus, in the $`{}_{}{}^{3}S_{1}^{}^3D_1`$ channel the renormalized solution at NLO in the Weinberg counting is divergent while LO and NNLO are convergent. In conclusion, the present analysis shows in a somewhat complementary manner as done in Sect. VIII.2 that indeed the NLO is problematic, at least non-perturbatively. In Sect. IX we will see that the problem is not solved if the NLO contribution is computed within a perturbative framework using the exact OPE-distorted wave basis as a lowest order approximation.
### VIII.4 The role of relativity and the $`\mathrm{\Delta }`$ resonance in the renormalization problem
The requirement of renormalizability may be regarded as a radical step, and renormalized LO calculations demand violating dimensional power counting on the counterterms Nogga et al. (2005) in non central waves such as $`{}_{}{}^{3}P_{0}^{}`$, due to an attractive $`1/r^3`$ singularity. To reach a finite limit the authors of Ref. Nogga et al. (2005) propose to promote counterterms which in Weinbergโs power counting are of higher order. However, in their proposal it is intriguing that they choose to promote just one counterterm in coupled channels, while they could have used a coupled channel counterterm, i.e. three counterterms in total. In the boundary condition approach we know from the start how many independent parameters must be exactly taken to reach a finite and unique limit, the reference to power counting is only specified at the level of the potential. Note that the power counting in the potential fixes its short attractive-repulsive singular character, and this is the origin of the conflict of assuming an a priori power counting for the counterterms. Finiteness requires that some forbidden counterterms must be allowed (promoted) Nogga et al. (2005) but also that some allowed counterterms must be forbidden (demoted). In such a framework, our NLO calculations in the $`{}_{}{}^{3}S_{1}^{}^3D_1`$ channel lead to finite but nonsensical results due to the repulsive-repulsive $`1/r^5`$ singularity (See Sect. VIII). On the other hand, if one fixes as required by the power counting the scattering length, the limit does not exist because one is driven to the exponentially diverging solution at the origin (For instance, Eq. (67) gives $`r_0\mathrm{}`$). How then can we reconcile finiteness with fixing of the parameters ?. As we pointed out already, the singular short distance behaviour of the chiral potential is in fact a long distance feature which changes dramatically when changing the long distance physics. Actually, one may reverse the argument and use renormalizability as a selective criterium for admissible long distance potentials. In the following we want to provide at least two possible scenarios how this might happen, i.e. , how modifying the potential at long distances by introducing physically relevant information the short distance behaviour of the potential changes.
In the first place, the chiral potential, Eq. (1), was derived in the heavy baryon expansion. The short distance character may change when not taking such a limit since the combination $`Mr`$ does make the order of limits ambiguous. A proper treatment of relativistic effects requires inclusion of antinucleons in loops, and a satisfactory EFT treatment of relativistic effects remains a challenging open problem because the standard crossing vs. unitarity non-perturbative divorce. On top of this one should use a satisfactory relativistic two body equation, which necessarily makes the problem fully non-local in coordinate space. Nevertheless, there exist โrelativisticโ potentials where some of the terms of higher power in 1/M than the TPE obtained in heavy-baryon ChPT are kept Higa and Robilotta (2003); Higa et al. (2004); Higa (2004) which have $`1/r^7`$ Van der Waals short distance behaviour with attractive-repulsive eigen potentials Higa (2005) meaning that as in the OPE case one has one free parameter. A calculation using these incomplete โrelativisticโ potentials will be presented elsewhere Higa et al. (2006).
A second scenario is related to the role played by the $`\mathrm{\Delta }`$ resonance <sup>15</sup><sup>15</sup>15We thank D. Phillips for drawing our attention to this point. not included in the present analysis. As pointed out in Ref. Kaiser et al. (1998), the $`\mathrm{\Delta }`$ provides the bulk of the chiral constants, yielding $`c_3=2c_4=g_A^2/2\mathrm{\Delta }`$, with $`\mathrm{\Delta }=293\mathrm{M}\mathrm{e}\mathrm{V}`$ the nucleon-delta mass splitting, yielding $`c_3=2.7\mathrm{GeV}^1`$ and $`c_4=1.35\mathrm{GeV}^1`$. The difference to the parameters of Ref. Entem and Machleidt (2002b) may be due to some other resonances. On the other hand, in terms of scales one has $`\mathrm{\Delta }2m_\pi `$, which might be regarded as a small parameter. This obviously does not mean that $`\mathrm{\Delta }`$ vanishes in the chiral limit. In the standard chiral counting of the potential, Eq. (1), the combinations $`\overline{c}_1=Mc_1`$, $`\overline{c}_3=Mc_3`$ and $`\overline{c}_4=Mc_4`$ are considered to be zeroth order, but according to the previous argument they could be regarded to be enhanced by one negative power. Thus the nominally NNLO terms containing $`c_3`$ and $`c_4`$ might become NLO contributions, and hence changing the repulsive-repulsive $`1/r^5`$ singularity into an attractive-attractive $`1/r^6`$ one. On the other hand, the $`c_3`$ and $`c_4`$ contributions of the standard NNLO dominate the short distance Van der Waals contributions. Actually, much of the NNLO potential is built from these terms all over the range. According to this reasoning our NNLO calculation may be closer to a NLO one where the $`N\mathrm{\Delta }`$ splitting is regarded as small parameter. In fact, taking NLO+$`\mathrm{\Delta }`$ with $`c_3=2c_4=g_A^2/2\mathrm{\Delta }`$ and $`\eta =0.0256`$ one gets $`A_S=0.8869\mathrm{fm}^{1/2}`$, $`Q_D=0.2762\mathrm{fm}^2`$, $`r_m=1.9726\mathrm{fm}`$ and $`P_d=0.06`$ in overall agreement with Table 3. It would be rather interesting to look for further consequences of this $`\mathrm{\Delta }`$-counting at higher orders. The importance of the $`\mathrm{\Delta }`$ in the NN problem has been stressed in several works already on power counting grounds Van Kolck (1993); Ordonez et al. (1996); Epelbaum et al. (2000); Pandharipande et al. (2005) but the crucial role played on the renormalization problem, i.e., the fact that the cut-off can be completely removed has not been recognized. Our discussion suggests that the momentum space cut-off could also be confortably removed in this $`\mathrm{\Delta }`$counting, unlike the delta-less NLO.
The two possible scenarios outlined above do not prove that the requirement of renormalizability is necessarily right, but suggest that looking into the short distance singular behaviour of long distance chiral potentials together with the mathematical requirement of finiteness may provide a significant physical insight into the NN problem. In the language of Ref. Nogga et al. (2005) where promotion of counterterms on the basis of the renormalizability requirement has been stressed, we are perhaps led also to the demotion of counterterms (like for relativistic potentials), or alternatively the promotion of terms in the potential (like in the $`\mathrm{\Delta }`$ counting described above).
### VIII.5 Van der Waals forces, the molecular analogy and the chiral quark model
The previous arguments show that it is possible to change the attractive/repulsive character of the potential at short distances by organizing the calculation of the potential in a different manner, but does not give a clue on why this actually happens. Remarkably, the analogy with atomic neutral systems subjected to Van der Waals forces illustrated in Sect. VII goes further and provides valuable insight into the problem. In low energy molecular physics where one works in a Born-Oppenheimer approximation, all atomic constituents, electrons and nuclei interact through the Coulomb force arising from one photon exchange. At long distances between distant electrons the potential is a dipole-dipole interaction
$`V_{\mathrm{dip}}(R)=e^2{\displaystyle \underset{A,B}{}}\left[{\displaystyle \frac{\stackrel{}{r}_A\stackrel{}{r}_B}{R^3}}3{\displaystyle \frac{(\stackrel{}{r}_A\stackrel{}{R})(\stackrel{}{r}_B\stackrel{}{R})}{R^5}}\right]`$ (93)
where the sum runs over electrons belonging to different atoms. In second order perturbation theory the atom-atom energy at a separation distance $`R`$ reads,
$`V_{AA}`$ $`=`$ $`AA|V_{\mathrm{dip}}|AA`$ (94)
$`+`$ $`{\displaystyle \underset{AAA^{}A^{}}{}}{\displaystyle \frac{|AA|V_{\mathrm{dip}}|A^{}A^{}|^2}{E_{AA}E_{A^{}A^{}}}}+\mathrm{}`$
where $`|AA`$ and $`|A^{}A^{}`$ is the electron wave function corresponding to a pair of separated clusters in their atomic ground state and excited states respectively. The first order contribution vanishes for atoms with no permanent dipole moment. The mutual electric polarization causes the Van der Waals interaction between the two atoms, $`C_6/R^6`$ and because it is second order perturbation theory it is obvious that the $`C_6`$ contribution to the potential will always be attractive. However, it is not clear that higher order terms would always be attractive. It is remarkable that the theorem of Thirring and Lieb Lieb and Thirring (1986) establishes that the Coulomb force between constituents implies all terms in the expansion being attractive, without appealing to the dipole-approximation. Thus, according to this result the long distance force will always be singular and attractive at short distances, and that is exactly what one needs. In such a situation making a long distance expansion of the potential, $`U=R_6^4/r^6R_8^6/r^8+\mathrm{}`$ and computing the scattering phase shifts by fixing always the same scattering length, along the lines pursued in this paper, makes much sense. Moreover, one expects the results for the phase shifts to be convergent if there is scale separation between the corresponding Van der Waals radii $`R_6R_8\mathrm{}`$. Our experience with several atomic systems confirms these expectations Calle Cordรณn et al. (2006).
The argument in the NN system is a straightforward generalization of the molecular system above. It is well known that there are no colour hidden states between colour neutral systems, so that at long distances one may assume only exchange of colourless objects. The longest range object will be the pion, and the mutual (chiral) polarizability will cause attraction between the nucleons, exactly in the same way as for atom-atom interactions. If we use as an example the chiral quark model, assume for simplicity non-relativistic constituent quarks one obtains the OPE for quarks. To second order perturbation theory we get the NN potential in the Born-Oppenheimer approximation
$`V_{NN}`$ $`=`$ $`NN|V_{\mathrm{OPE}}|NN`$ (95)
$`+`$ $`{\displaystyle \underset{HHNN}{}}{\displaystyle \frac{|NN|V_{\mathrm{OPE}}|HH|^2}{E_{NN}E_{HH}}}+\mathrm{}`$
where $`V_{NN}`$ represents the potential in the $`NN`$ operator basis. This yields exactly when $`HH=N\mathrm{\Delta }`$ the results found in Ref. Van Kolck (1993); Ordonez et al. (1996); Kaiser et al. (1998) and naturally explains why the contribution from one $`\mathrm{\Delta }`$ intermediate state is attractive at short distances. Although this analogy with molecular systems is very suggestive, the generalization to all orders along the lines of the Lieb-Thirring theorem within a QCD context remains at present an optimistic speculation.
## IX Renormalized Perturbation Theory versus non-integer power counting
### IX.1 Perturbations on boundary conditions
In all our calculations we have taken a long distance potential calculated perturbatively, and scattering amplitudes have been computed non-perturbatively by fully iterating a potential computed in perturbation theory, as initially suggested by Weinberg Weinberg (1990). We will call this form of solution non-perturbative for brevity. This requires a non-perturbative treatment of the renormalization problem, which naturally implies that the short distance renormalization conditions (or counterterms) are determined by the most singular contribution of the long distance potential at the origin. For the NN chiral potential it turns out that the higher the order the more singular the potential. As a consequence we have seen in Sect. VIII.3 that for instance Weinberg counting at NLO in the $`{}_{}{}^{3}S_{1}^{}^3D_1`$ channel is incompatible with renormalization and finiteness due to the short distance repulsive character of the NLO potential. Although naively this looks counterintuitive, it is important to realize that there are also situations, like LO and NNLO, where the regularity condition of the wave function conspires against the singularity so that the net effect is well behaved in the scattering amplitudes and deuteron properties. Our results in Sect. IV and V, suggest that the pattern obtained when comparing LO and NNLO looks quite converging numerically, although there appears to be no way of making an a priori estimate of the corrections.
Perturbative treatments might circumvent this difficulty since they have the indubitable benefit of allowing an a priori estimate of the systematic error via dimensional power counting. This causes no problem in the calculation of the long distance potential. However, we anticipate already that singular potentials are indeed singular perturbations, and power counting may not work as one naively expects for the full amplitudes. Kaplan, Savage and Wise Kaplan et al. (1998, 1999) suggested such a perturbative scheme some years ago, where the lowest order approximation was a contact theory and OPE and higher order corrections could be computed in perturbation theory. This is equivalent to consider $`mM/f^2`$ to be first order and $`m^2/f^2`$ second order, so that a calculation involving the chiral constants would be N<sup>3</sup>LO in that counting. Unfortunately, the expansion turned out to be non-converging at NNLO Fleming et al. (2000). In our coordinate space formulation, this approach corresponds to assume for the $`S`$ waves a boundary condition fixing the scattering length $`\alpha _0`$ Pavon Valderrama and Ruiz Arriola (2005a) (See Appendix A of that work) and making long distance potential perturbations. In our previous work we verified that perturbation theory could only account for a contribution to the deuteron and $`{}_{}{}^{3}S_{1}^{}^3D_1`$ scattering observables at first order. Unfortunately, the second order was divergent, while non-perturbatively, i.e., exactly solving the Schrรถdinger equation for the OPE potential, the results were not only finite but also numerically quite close to experiment. This deserves some explanation. In Fig. 15 we show the results in the deuteron channel when we scale the OPE potential $`U_{\mathrm{OPE}}\lambda U_{\mathrm{OPE}}`$ for the s-wave function normalization $`A_S(\lambda )`$ and the effective range $`r_0(\lambda )`$ as a function of the scaling parameter $`\lambda `$ by keeping the deuteron binding energy fixed to its experimental value. The non-perturbative result is compared to the first order perturbation theory used in Ref. Pavon Valderrama and Ruiz Arriola (2005a). Clearly, perturbation theory fails even for weak coupling. The experimental value for $`r_0`$ could be obtained by adding a counterterm $`C_2`$ as done by Kaplan, Savage and Wise Kaplan et al. (1998, 1999). A non-vanishing $`C_2`$ not only violates the orthogonality of the zero energy and deuteron wave functions for a long distance local potential but also introduces a new parameter, reducing the predictive power. Moreover, the non-perturbative inclusion of this $`C_2`$ counterterm with the OPE potential yields divergent results (see the discussion in Sect. VIII.3). In fact, much of the strength of $`C_2`$ is naturally provided by the short distance $`1/r^3`$ singularity of the OPE potential.
### IX.2 Perturbations on the OPE potential
Recently, Nogga,Timmermans and van Kolck (NTvK) have suggested Nogga et al. (2005) treating the OPE effects non-perturbatively, i.e. to all orders, while TPE and higher as well as $`\mathrm{\Delta }`$ contributions should be computed in perturbation theory (see also Refs. Griesshammer (2005); Birse (2005a, b) for related ideas). In this section, we analyze such a proposal disregrading the Delta. Our main conclusions will not change although numbers could be modified. In such a situation the perturbative expansion is equivalent to consider the $`mM/(4\pi f^2)`$ to be zeroth order while $`m^2/(4\pi f)^2`$ is taken to be second order and $`m/M`$ is first order, so that the potential can written as
$`U^{(0)}`$ $`=`$ $`U_{1\pi }^{(0)},`$
$`U^{(2)}`$ $`=`$ $`U_{1\pi }^{(2)}+U_{2\pi }^{(2)},`$
$`U^{(3)}`$ $`=`$ $`U_{1\pi }^{(3)}+U_{2\pi }^{(3)}.`$ (96)
Non-perturbatively $`U_{1\pi }=U_{1\pi }^{(0)}+U_{2\pi }^{(0)}+U_{3\pi }^{(0)}+\mathrm{}`$ amounts to take $`g_{\pi NN}=13.1`$ in the OPE piece, hence accounting for the Goldberger-Treiman discrepancy. In perturbation theory, we must take $`U_{1\pi }^{(0)}`$ with $`g_A=1.26`$ and include OPE corrections to higher order. Note that the missing first order implies substantial simplifications in the perturbative treatment. Indeed NNLO can be done within first order perturbation theory since going to second order perturbation theory considers $`U_{}^{(2)}{}_{}{}^{2}`$ which is N<sup>3</sup>LO. In the remaining of this section we analyze some aspects of such proposal by scaling the strength of the perturbation and show that the appearance of non-analytical behaviour is intrinsic to singular potentials, yielding to perturbative divergences. As we will see, finite perturbative calculations including chiral TPE to NNLO would require 4 counterterms for the singlet $`{}_{}{}^{1}S_{0}^{}`$ and 6 counterterms for the triplet $`{}_{}{}^{3}S_{1}^{}^3D_1`$ channel. Our non-perturbative results, i.e., fully iterated NNLO potentials, are based on just 1 and 3 counterterms respectively.
### IX.3 Singlet $`{}_{}{}^{1}S_{0}^{}`$ channel in distorted OPE waves
Let us examine first the $`{}_{}{}^{1}S_{0}^{}`$ channel, and consider the effect of the NLO and NNLO TPE potentials on top of the LO OPE potential in long distance perturbation theory. The fact that they are taken second and third order respectively means that the effect will be additive at NNLO in the scattering properties. For the total potential in Eq. (96) we write the wave function as
$`u_k(r)=u_k^{(0)}(r)+u_k^{(2)}(r)+u_k^{(3)}(r)+\mathrm{}\mathrm{}`$ (97)
and the phase shift becomes
$`\delta _0=\delta _0^{(0)}+\delta _0^{(2)}+\delta _0^{(3)}+\mathrm{}`$ (98)
At LO the scattering length $`\alpha _0^{(0)}`$ is a free parameter which we fix to the physical value, $`\alpha _0^{(0)}=\alpha _0`$. As we did in our non-perturbative treatment in Sect. VIII.1, we will keep the scattering length fixed to its experimental value at any order of the approximation, so that differences may be only attributable to the potential. In the normalization of Eq. (44) the correction to the phase shift is just given by
$`\delta _0^{(2)}=k\mathrm{sin}^2\delta _0^{(0)}{\displaystyle _{r_c}^{\mathrm{}}}U^{(2)}(r)u_k^{(0)}(r)^2๐r,`$ (99)
and a similar expression for $`\delta _0^{(3)}`$ which can be deduced by the standard Lagrangeโs identity. Here, a short distance cut-off $`r_c`$ has been assumed because at short distances the NLO potential diverges as $`U_{\mathrm{NLO}}1/r^5`$. The previous formula yields a change also in the scattering length, so that we may eliminate the cut-off radius by subtracting off the zero energy contribution by fixing $`\alpha _0^{(2)}=0`$. It is convenient to recast the result in the form of an effective range expansion in the OPE distorted wave basis,
$`k\mathrm{cot}\delta _0+{\displaystyle \frac{1}{\alpha _0}}`$ $`=`$ $`k\mathrm{cot}\delta _0^{(0)}+{\displaystyle \frac{1}{\alpha _0^{(0)}}}`$
$`+`$ $`{\displaystyle _{r_c}^{\mathrm{}}}๐rU^{(2)}(r)\left[u_k^{(0)}(r)^2u_0^{(0)}(r)^2\right]`$
$`+`$ $`{\displaystyle _{r_c}^{\mathrm{}}}๐rU^{(3)}(r)\left[u_k^{(0)}(r)^2u_0^{(0)}(r)^2\right],`$
which guarantees $`\alpha _0^{(2)}=\alpha _0^{(3)}=0`$, due to the one subtraction. If we expand in powers of the energy the LO wave function we get
$`u_k^{(0)}(r)=u_0^{(0)}(r)+k^2u_2^{(0)}(r)+k^4u_4^{(0)}(r)+\mathrm{}.`$ (101)
where
$`u_0^{(0)}{}_{}{}^{\prime \prime }(r)+U(r)u_0^{(0)}(r)`$ $`=`$ $`0`$
$`u_2^{(0)}{}_{}{}^{\prime \prime }(r)+U(r)u_2^{(0)}(r)`$ $`=`$ $`u_0^{(0)}(r)`$
$`u_4^{(0)}{}_{}{}^{\prime \prime }(r)+U(r)u_4^{(0)}(r)`$ $`=`$ $`u_2^{(0)}(r)`$
$`\mathrm{}`$ (102)
These equations can be solved recursively. Thus, the NLO correction to the effective range is given by
$`r_0^{(2)}=4{\displaystyle _{r_c}^{\mathrm{}}}U^{(2)}(r)u_2^{(0)}(r)u_0^{(0)}(r)๐r.`$ (103)
To estimate the short distance contribution we use the OPE exchange potential in the form $`U_{\mathrm{LO}}=e^{mr}/(R_sr)`$, with $`R_s=16f^2\pi /g^2m^2M`$ the characteristic length $`{}_{}{}^{1}S_{0}^{}`$-channel scale. Note that the OPE potential in the $`{}_{}{}^{1}S_{0}^{}`$ channel is Coulomb like at short distances for which the complete regular plus irregular solution is known. One could then use a short distance expansion of the general analytical Coulomb solution. This facilitates guessing the solution at short distances for the zeroth energy wave function. The higher energy wave functions can be computed straightforwardly, yielding for $`r0`$
$`u_0^{(0)}(r)`$ $``$ $`c_0\left[1+mr{\displaystyle \frac{3r}{2R_s}}{\displaystyle \frac{r}{R_s}}\mathrm{log}\left({\displaystyle \frac{r}{R_s}}\right)\right]+c_1r`$
$`u_2^{(0)}(r)`$ $``$ $`c_0rR_s+๐ช(r^3)`$
$`u_4^{(0)}(r)`$ $``$ $`{\displaystyle \frac{1}{3!}}c_0r^3R_s+๐ช(r^5)`$
$`u_6^{(0)}(r)`$ $``$ $`{\displaystyle \frac{1}{5!}}c_0r^5R_s+๐ช(r^7)`$
$`u_8^{(0)}(r)`$ $``$ $`{\displaystyle \frac{1}{7!}}c_0r^7R_s+๐ช(r^9)`$ (104)
as can be readily checked by solving the Schrรถdinger equation in powers of energy, Eq. (102). The coefficients $`c_1`$ and $`c_0`$ correspond to the linearly independent regular and irregular solutions respectively and are determined by matching to the integrated in asymptotic condition $`u_0^{(0)}(r)1r/\alpha _0`$ at large distances and at zero energy. Obviously, the irregular solution contributes, $`c_00`$, because $`\alpha _0`$ is taken to be independent of the potential, and hence terms proportional to the coefficient $`c_1`$ are subleading. Thus, we get
$`r_0^{(2)}4{\displaystyle _{r_c}}๐r{\displaystyle \frac{MC_5}{r^5}}(c_0^2R_s)r.`$ (105)
Thus, we conclude that the first order perturbative result is badly divergent. This is very puzzling since the non-perturbative calculation in Sect. VIII.1 yields a finite number (see Eq. (89)), and suggests non-analytical dependence on the coupling constant. To enlighten the situation, let us scale the NLO potential by a factor $`\lambda `$, $`U_{\mathrm{NLO}}\lambda U_{\mathrm{NLO}}`$, and compute non-perturbatively the effective range as a function of the scaling parameters, $`r_0(\lambda )`$, with the obvious conditions $`r_0(0)=r_0^{(0)}`$ and $`r_0(1)=r_0^{\mathrm{NLO}}`$. The result is presented in Fig. 16. The infinite slope at the origin can be clearly seen. Numerically, we find that for small $`\lambda 0.1`$, the correction to the effective range behaves as $`r_0r_0^{(0)}\sqrt{\lambda }`$, whereas for $`\lambda 1`$, it behaves as $`r_0r_0^{(0)}\lambda ^{\frac{1}{3}}`$. This fractional power counting $`\lambda ^\alpha `$, with $`0<\alpha <1`$ is evident from the universal low energy theorem, Eq. (39), for a potential with a single scale, $`U(r)=F(r/R)/R^2`$, and appears when the short distance regulator is removed <sup>16</sup><sup>16</sup>16In Ref. Beane et al. (2001) a potential well was used as a short distance regulator which was not removed, and hence the non-analiticity was not seen.. An explicit example is provided by Eq. (82) when a pure Van der Waals potential acts as a perturbation to a boundary condition ( a contact theory with only $`\alpha _0`$), since the strength of the potential is $`\lambda MC_6=R^4`$, but $`r_0`$ contains $`R\lambda ^{1/4}`$, $`R^2\lambda ^{1/2}`$ and $`R^3\lambda ^{3/4}`$. It would be interesting to predict a priori this non-perturbative non-integer power counting analytically for potentials with multiple scales as we have done here numerically Calle Cordรณn et al. (2006) <sup>17</sup><sup>17</sup>17Fractional power counting has also been reported to occur also in the EFT analysis of the three body problem for the pionless theory Griesshammer (2005); Birse (2005b).
Obviously, to prevent the perturbative divergence one could subtract an energy dependent contribution, and provide the effective range as an input parameter <sup>18</sup><sup>18</sup>18This is equivalent to use a short distance energy dependent boundary condition on the solution and hence to violate the orthogonality conditions discussed in Sect. III.. Then one would get
$`k\mathrm{cot}\delta `$ $`=`$ $`k\mathrm{cot}\delta ^{(0)}+{\displaystyle \frac{1}{2}}\left(r_0r_0^{(0)}\right)k^2`$ (106)
$`+`$ $`{\displaystyle _{r_c}^{\mathrm{}}}dr[U^{(2)}(r)+U^{(3)}(r)]\times `$
$`\left[u_k^{(0)}(r)^2u_0^{(0)}(r)^22k^2u_0^{(0)}(r)u_2^{(0)}(r)^2\right]`$
Note that this equation requires assuming $`r_0r_0^{(0)}=๐ช(\lambda )`$, while non-perturbatively we find $`r_0r_0^{(0)}=๐ช(\lambda ^{\frac{1}{2}})`$. Now the NLO and NNLO corrections to the $`v_2`$ parameter would come as a prediction,
$`v_2^{(2)}`$ $`+`$ $`v_2^{(3)}={\displaystyle _{r_c}^{\mathrm{}}}dr[U^{(2)}(r)+U^{(3)}(r)]\times `$ (107)
$`\left[2u_4^{(0)}(r)u_0^{(0)}(r)+u_2^{(0)}(r)^2\right],`$
which is also divergent since the leading behaviour of the integrand is $`1/r^3`$ at NLO and $`1/r^4`$ at NNLO for small $`r`$, see Eq. (104). Thus, a further subtraction would be needed, predicting the correction to $`v_3`$,
$`v_3^{(2)}`$ $`+`$ $`v_3^{(3)}={\displaystyle _{r_c}^{\mathrm{}}}dr[U^{(2)}(r)+U^{(3)}(r)]\times `$
$`\left[2u_6^{(0)}(r)u_0^{(0)}(r)+2u_2^{(0)}(r)u_4^{(0)}(r)\right],`$
which is logarithmically divergent due to Eq. (104). Finally, if a fourth subtraction is implemented a convergent prediction is obtained for $`v_4`$ at NLO and NNLO,
$`v_4^{(2)}`$ $`+`$ $`v_4^{(3)}={\displaystyle _{r_c}^{\mathrm{}}}dr[U^{(2)}(r)+U^{(3)}(r)]\times `$
$`\left[2u_8^{(0)}(r)u_0^{(0)}(r)+2u_2^{(0)}(r)u_4^{(0)}(r)+u_4^{(0)}(r)^2\right].`$
These four subtractions, needed to make a renormalized perturbative prediction of the $`{}_{}{}^{1}S_{0}^{}`$ phase shift at NNLO , actually correspond to having 4 counterterms, i.e. fixing $`\alpha _0,r_0,v_2`$ and $`v_3`$. This result disagrees with the standard Weinberg counting (two counterterms at NLO and NNLO in the $`{}_{}{}^{1}S_{0}^{}`$ channel). Moreover, besides the loss of predictive power as compared to the non-perturbative result where only one counterterm is needed the deduced renormalized value for $`v_4`$ is worsened in perturbation theory, since $`v_4^{(0)}=50.74\mathrm{fm}^7`$, $`v_4^{(2)}=10.45\mathrm{fm}^7`$ and $`v_4^{(3)}=2.88\mathrm{fm}^7`$. The situation is summarized in Table 7 where we show our numerical results obtained in perturbation theory as explained above and they are compared to the NijmII and Reid93 potential model calculations (See e.g. Ref. Pavon Valderrama and Ruiz Arriola (2004c)). Although these are not directly experimental data it is noteworthy that they differ by a few percent while the perturbative calculation is about a factor 3 larger. The integrals for $`v_4`$ are rather well converging and the matching between the numerical solution and the short distance solutions, Eq. (104), is quite stable in the region around $`r0.1\mathrm{fm}`$.
Thus, in this particular example of the $`{}_{}{}^{1}S_{0}^{}`$ channel one sees that our non-perturbative approach based on the choice of the regular solutions at the origin predicts the phase shift and hence all low energy parameters from $`\alpha _0`$ and the potential as displayed in Table 6. A perturbative treatment of the amplitude based on OPE distorted waves requires to fix $`\alpha _0,r_0,v_2`$ and $`v_3`$ at NNLO. The phenomenological success and converging pattern observed when the potential is considered at LO,NLO and NNLO is solved non-perturbatively is very encouraging. The price to pay is to face non-analytical behaviour which implies a non-integer power counting. The trend observed here can be generalized to other channels. A more thorough discussion of this issue will be presented elsewhere Calle Cordรณn et al. (2006).
### IX.4 Triplet $`{}_{}{}^{3}S_{1}^{}^3D_1`$ channel in distorted OPE waves
We turn now to the triplet $`{}_{}{}^{3}S_{1}^{}^3D_1`$ channel. The reasoning is a straightforward, although tedious, coupled channel generalization of the $`{}_{}{}^{1}S_{0}^{}`$ case, with the additional feature that the short distance behaviour is dominated by a $`1/r^3`$ singularity (instead of $`1/r`$), and so the short distance behaviour is different. It is convenient to introduce the potential matrix as
$`๐(r)`$ $`=`$ $`\left(\begin{array}{cc}U_{{}_{}{}^{3}S_{1}^{}}(r)& U_{E_1}(r)\\ U_{E_1}(r)& U_{{}_{}{}^{3}D_{1}^{}}(r)\end{array}\right),`$ (110)
and the matrix wave function,
$`๐ฎ_k(r)`$ $`=`$ $`๐\left(\begin{array}{cc}u_{k,\alpha }(r)& u_{k,\beta }(r)\\ w_{k,\alpha }(r)& w_{k,\beta }(r)\end{array}\right),`$ (111)
with $`๐`$ a constant energy dependent matrix, subject to a slightly different normalization than Eq. (LABEL:eq:phase\_triplet),
$`๐ฎ_k(r){\displaystyle \frac{1}{k}}\widehat{๐ฃ}(kr)๐^1\widehat{๐}\widehat{๐ฒ}(kr)๐`$ (112)
Here, $`\widehat{๐}`$ is the effective range matrix defined by its relation to the unitary $`๐`$-matrix,
$`\mathrm{๐๐๐}^1=\left(\widehat{๐}+\mathrm{i}k๐^2\right)\left(\widehat{๐}\mathrm{i}k๐^2\right)^1,`$ (113)
and $`๐=\mathrm{diag}(1,k^2)`$. The reduced Bessel functions matrices are given by $`\widehat{๐ฃ}=\mathrm{diag}(\widehat{j}_0,\widehat{j}_2)`$ and $`\widehat{๐ฒ}=\mathrm{diag}(\widehat{y}_0,\widehat{y}_2)`$ with $`\widehat{j}_l(x)=xj_l(x)`$ and $`\widehat{y}_l(x)=xy_l(x)`$. At low energies, one has the effective range expansion(see e.g. Pavon Valderrama and Ruiz Arriola (2004c) and references therein),
$`\widehat{๐}=(๐)^1+{\displaystyle \frac{1}{2}}๐ซk^2+๐ฏk^4+\mathrm{}`$ (114)
Here, we have introduced the scattering length matrix,
$`๐`$ $`=`$ $`\left(\begin{array}{cc}\alpha _0& \alpha _{02}\\ \alpha _{02}& \alpha _2\end{array}\right),`$ (115)
the effective range matrix
$`๐ซ`$ $`=`$ $`\left(\begin{array}{cc}r_0& r_{02}\\ r_{02}& r_2\end{array}\right),`$ (116)
and so on. These parameters have been determined in Pavon Valderrama and Ruiz Arriola (2004c) from the potentials of Ref. Stoks et al. (1994). Proceeding similarly as in the one channel case, one gets, after one subtraction at zero energy the effective range function in perturbation theory,
$`\widehat{๐}`$ $`+`$ $`(๐)^1=\widehat{๐}^{(0)}+(๐^{(0)})^1`$ (117)
$`+`$ $`{\displaystyle _{r_c}^{\mathrm{}}}๐r\left[๐ฎ_k^{(0)}๐^{(2)}๐ฎ_k^{(0)}๐ฎ_0^{(0)}๐^{(2)}๐ฎ_0^{(0)}\right]`$
The condition $`\alpha _0^{(0)}=\alpha _0`$ must be imposed, since $`\alpha _{02}^{(0)}`$ and $`\alpha _2^{(0)}`$ are predicted from $`\alpha _0^{(0)}`$ (at LO one only needs one counterterm). This formula implies that one introduces two new conditions to fix now $`\alpha _{02}`$ and $`\alpha _2`$ to their experimental value.Along similar lines as done before, we analyze the finiteness of the previous expression by computing the effective range matrix. To this end we expand the coupled channel wave function in powers of momentum
$`๐ฎ_k^{(0)}(r)=๐ฎ_0^{(0)}(r)+k^2๐ฎ_2^{(0)}(r)+k^4๐ฎ_4^{(0)}(r)+\mathrm{}`$ (118)
to get
$`๐ซ^{(2)}={\displaystyle _{r_c}^{\mathrm{}}}๐r\left[๐ฎ_2^{(0)}๐^{(2)}๐ฎ_0^{(0)}+๐ฎ_0^{(0)}๐^{(2)}๐ฎ_2^{(0)}\right]`$ (119)
The LO OPE short distance behaviour of the triplet wave functions has been worked out in our previous work Pavon Valderrama and Ruiz Arriola (2005a). It is convenient to define the triplet length scale
$`R_t={\displaystyle \frac{3g_A^2M}{32\pi f_\pi ^2}}`$ (120)
which value $`R_t=1.07764\mathrm{fm}`$. One has the general structure
$`u(r)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{3}}}\left({\displaystyle \frac{r}{R_t}}\right)^{3/4}[C_{1R}f_{1R}(r)e^{+4\sqrt{2}\sqrt{\frac{R_t}{r}}}`$
$``$ $`C_{2R}f_{2R}(r)e^{4\sqrt{2}\sqrt{\frac{R_t}{r}}}+\sqrt{2}C_{1A}f_{1A}(r)e^{4i\sqrt{\frac{R_t}{r}}}`$
$`+`$ $`\sqrt{2}C_{2A}f_{2A}(r)e^{4i\sqrt{\frac{R_t}{r}}}]`$
$`w(r)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{3}}}\left({\displaystyle \frac{r}{R_t}}\right)^{3/4}[\sqrt{2}C_{1R}g_{1R}(r)e^{+4\sqrt{2}\sqrt{\frac{R_t}{r}}}`$
$`+`$ $`\sqrt{2}C_{2R}g_{2R}(r)e^{4\sqrt{2}\sqrt{\frac{R_t}{r}}}+C_{1A}g_{1A}(r)e^{4i\sqrt{\frac{R_t}{r}}}`$
$`+`$ $`C_{2A}g_{2A}(r)e^{4i\sqrt{\frac{R_t}{r}}}]`$
where the constants $`C_{1R}`$, $`C_{2R}`$, $`C_{1A}`$ and $`C_{2A}`$ depend on the energy and the OPE potential parameters. The regular solution is selected when one takes $`C_{1R}=0`$. The functions appearing in this formula are of the form
$`f(r)`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}a_n\left({\displaystyle \frac{r}{R_t}}\right)^{n/2}`$
$`g(r)`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}b_n\left({\displaystyle \frac{r}{R_t}}\right)^{n/2}`$
For the present calculation we only need the power behaviour (see Appendix B of Ref. Pavon Valderrama and Ruiz Arriola (2005a))
$`๐ฎ_0^{(0)}(r)`$ $``$ $`r^{3/4}`$
$`๐ฎ_2^{(0)}(r)`$ $``$ $`r^{3/4+5/2}`$
$`๐ฎ_4^{(0)}(r)`$ $``$ $`r^{3/4+5}`$ (123)
which shows that, again, the first order correction to the effective range matrix is logarithmically divergent because the NLO potential diverges as $`1/r^5`$ and $`๐^{(2)}๐ฎ_0๐ฎ_21/r`$ <sup>19</sup><sup>19</sup>19There is a subtlety here. The terms containing the regular exponential at the origin are convergent, regardless on the power of $`r`$ in the denominator. Naively logarithmically divergent integrals would become convergent when combined with oscillating functions. However, these functions appear squared so that the logarithmic divergence prevails.. As previously, the situation could be amended by adding 3 new counterterms to fix the effective range matrix $`๐ซ`$ and then $`๐ฏ`$ would come as a prediction. So, at NLO in perturbation theory one needs a total of 6 counterterms to generate a coupled channel finite amplitude. When adding the NNLO contribution this number of counterterms remains the same since $`๐^{(3)}๐ฎ_0^{(0)}๐ฎ_4^{(0)}r^{1/2}`$ and $`๐^{(3)}[๐ฎ_2^{(0)}]^2r^{1/2}`$.
An illustration of non-analytical non-perturbative behaviour in the $`{}_{}{}^{3}S_{1}^{}^3D_1`$-channel can be looked up in Fig. 17. There, the behaviour of the s-wave function normalization $`A_S(\lambda )`$ in the deuteron and the effective range $`r_0(\lambda )`$ in the $`{}_{}{}^{3}S_{1}^{}^3D_1`$ channel when the TPE potential is scaled as $`U=U_{\mathrm{LO}}+\lambda (U_{\mathrm{NLO}}+U_{\mathrm{NNLO}})`$ and the deuteron binding energy, the asymptotic $`D/S`$-ration, $`\eta `$ and the s-wave scattering length, $`\alpha _0`$ are fixed to their experimental values.
### IX.5 The deuteron in distorted OPE waves
To conclude our analysis of perturbation theory we study now the deuteron bound state. According to Fig. 17 there appears some tiny non-analyticity for very small couplings in the asymptotic s-wave normalization $`A_S`$. Note that there is an apparent linear behaviour with the exception of the very small $`\lambda `$ region, making one suspect that the result might be obtained in perturbation theory. We will see below by an explicit perturbative calculation that this is not so. We have checked that this trend also occurs for other quantities such as the quadrupole moment, $`Q_d`$, the matter radius, $`r_m`$, and the D-state probability, $`P_D`$. Here, we show, as it has been done above for the scattering problem, that this can be traced to a first order divergent renormalized result.
We define the two component deuteron state,
$`๐ฎ_\gamma (r)=\left(\begin{array}{c}u_\gamma (r)\\ w_\gamma (r)\end{array}\right)`$ (124)
In perturbation theory, we expand the potential
$`๐(r)=๐^{(0)}(r)+๐^{(2)}(r)+๐^{(3)}(r)+\mathrm{}`$ (125)
and thus the deuteron wave function for fixed energy (or $`\gamma `$),
$`๐ฎ_\gamma (r)=๐ฎ_\gamma ^{(0)}(r)+๐ฎ_\gamma ^{(2)}(r)+๐ฎ_\gamma ^{(3)}(r)+\mathrm{}`$ (126)
where $`(u_\gamma ^{(0)}(r),w_\gamma ^{(0)}(r))`$ correspond to the lowest order solutions of the problem and $`(u_\gamma ^{(2)}(r),w_\gamma ^{(2)}(r))`$ and $`(u_\gamma ^{(3)}(r),w_\gamma ^{(3)}(r))`$ satisfy
$`๐ฎ_\gamma ^{(0)}(r)+\left[๐^{(0)}(r)+\gamma ^2\right]๐ฎ_\gamma ^{(0)}(r)`$ $`=`$ $`0`$
$`๐ฎ_\gamma ^{(2)}(r)+\left[๐^{(0)}(r)+\gamma ^2\right]๐ฎ_\gamma ^{(2)}(r)`$ $`=`$ $`๐^{(2)}(r)๐ฎ_\gamma ^{(0)}(r)`$
$`๐ฎ_\gamma ^{(3)}(r)+\left[๐^{(0)}(r)+\gamma ^2\right]๐ฎ_\gamma ^{(3)}(r)`$ $`=`$ $`๐^{(3)}(r)๐ฎ_\gamma ^{(0)}(r)`$
We look for normalized solutions, so that perturbatively
$`1`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}๐r๐ฎ_\gamma ^{(0)}(r)๐ฎ_\gamma ^{(0)}(r)`$
$`0`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}๐r\left(๐ฎ_\gamma ^{(2)}(r)๐ฎ_\gamma ^{(0)}(r)+๐ฎ_\gamma ^{(0)}(r)๐ฎ_\gamma ^{(2)}(r)\right)`$
$`0`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}๐r\left(๐ฎ_\gamma ^{(3)}(r)๐ฎ_\gamma ^{(0)}(r)+๐ฎ_\gamma ^{(0)}(r)๐ฎ_\gamma ^{(3)}(r)\right)`$
The zeroth order equation was solved in our previous work Pavon Valderrama and Ruiz Arriola (2005a), where it was shown that $`\gamma `$ was a free parameter, which means $`\gamma ^{(0)}=\gamma `$, and the regular solution at the origin was selected (see Eq. (LABEL:eq:LO\_short\_wf)) to ensure normalizability at the origin. We will keep always the same fixed value at any order of the approximation, so that $`\gamma ^{(2)}=\gamma ^{(3)}=0`$. To analyze the NLO and NNLO problem analytically we proceed by the variable coefficients method. The zeroth order equation is a homogenous linear system with four linearly independent solutions,
$`๐ฎ_i^{(0)}(r)=\left(\begin{array}{c}u_i(r)\\ w_i(r)\end{array}\right)i=1,2,3,4`$ (129)
The first order equation is an inhomogeneous linear system, which solution can be written as
$`u_\gamma ^{(2)}(r)+u_\gamma ^{(3)}(r)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{4}{}}}c_i(r)u_i(r)`$
$`w_\gamma ^{(2)}(r)+w_\gamma ^{(3)}(r)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{4}{}}}c_i(r)w_i(r)`$ (130)
The variable coefficients satisfy
$`{\displaystyle \underset{i=1}{\overset{4}{}}}c_i^{}(r)u_i(r)`$ $`=`$ $`0`$
$`{\displaystyle \underset{i=1}{\overset{4}{}}}c_i^{}(r)w_i(r)`$ $`=`$ $`0`$
$`{\displaystyle \underset{i=1}{\overset{4}{}}}c_i^{}(r)u_i^{}(r)`$ $`=`$ $`F_u(r)`$
$`{\displaystyle \underset{i=1}{\overset{4}{}}}c_i^{}(r)w_i^{}(r)`$ $`=`$ $`F_w(r)`$ (131)
where we have defined the driving term
$`๐
(r)`$ $`=`$ $`\left(\begin{array}{c}F_u(r)\\ F_w(r)\end{array}\right)`$ (132)
$`=`$ $`๐^{(2)}(r)๐ฎ_\gamma ^{(0)}(r)๐^{(3)}(r)๐ฎ_\gamma ^{(0)}(r)`$
which at short distances behaves as
$`๐
(r)r^{3/4}C_5r^5+r^{3/4}C_6r^6`$ (133)
whereas at large distances one has
$`๐
(r)e^{\gamma r}e^{mr}`$ (134)
To proceed further, we choose the following linearly independent solutions fulfilling the asymptotic boundary condition at infinity
$`u_1(r)`$ $``$ $`e^{\gamma r},`$
$`w_1(r)`$ $``$ $`0,`$
$`u_2(r)`$ $``$ $`0,`$
$`w_2(r)`$ $``$ $`e^{\gamma r}\left(1+{\displaystyle \frac{3}{\gamma r}}+{\displaystyle \frac{3}{(\gamma r)^2}}\right),`$
$`u_3(r)`$ $``$ $`e^{\gamma r},`$
$`w_3(r)`$ $``$ $`0,`$
$`u_4(r)`$ $``$ $`0,`$
$`w_4(r)`$ $``$ $`e^{\gamma r}\left(1{\displaystyle \frac{3}{\gamma r}}+{\displaystyle \frac{3}{(\gamma r)^2}}\right),`$
Any of these solutions has a short distance behaviour of the general form given in Eq. (LABEL:eq:LO\_short\_wf). So that, all these solutions are necessarily singular at the origin. Using Krammerโs rule the solutions to the linear differential system, Eq. (131), which are regular at infinity read
$`c_1(r)={\displaystyle \frac{1}{W}}{\displaystyle _0^r}๐r^{}\left|\begin{array}{cccc}0& u_2& u_3& u_4\\ 0& w_2& w_3& w_4\\ F_u& u_2^{}& u_3^{}& u_4^{}\\ F_w& w_2^{}& w_3^{}& w_4^{}\end{array}\right|`$ (140)
$`c_2(r)={\displaystyle \frac{1}{W}}{\displaystyle _0^r}๐r\left|\begin{array}{cccc}u_1& 0& u_3& u_4\\ w_1& 0& w_3& w_4\\ u_1^{}& F_u& u_3^{}& u_4^{}\\ w_1^{}& F_w^{}& w_3^{}& w_4^{}\end{array}\right|`$ (145)
$`c_3(r)={\displaystyle \frac{1}{W}}{\displaystyle _r^{\mathrm{}}}๐r^{}\left|\begin{array}{cccc}u_1& u_2& 0& u_4\\ w_1& w_2& 0& w_4\\ u_1^{}& u_2^{}& F_u& u_4^{}\\ w_1^{}& w_2^{}& F_w& w_4^{}\end{array}\right|`$ (150)
$`c_4(r)={\displaystyle \frac{1}{W}}{\displaystyle _r^{\mathrm{}}}๐r^{}\left|\begin{array}{cccc}u_1& u_2& u_3& 0\\ w_1& w_2& w_3& 0\\ u_1^{}& u_2^{}& u_3^{}& F_u\\ w_1^{}& w_2^{}& w_3^{}& F_w\end{array}\right|`$ (155)
where $`W`$ is the Wronskian
$`W=\left|\begin{array}{cccc}u_1& u_2& u_3& u_4\\ w_1& w_2& w_3& w_4\\ u_1^{}& u_2^{}& u_3^{}& u_4^{}\\ w_1^{}& w_2^{}& w_3^{}& w_4^{}\end{array}\right|=4\gamma ^2.`$ (160)
At asymptotically large distances we have
$`u^{(2)}(r)`$ $``$ $`c_S^{(2)}e^{\gamma r}`$
$`w^{(2)}(r)`$ $``$ $`c_D^{(2)}\eta ^{(0)}e^{\gamma r}\left(1+{\displaystyle \frac{3}{\gamma r}}+{\displaystyle \frac{3}{(\gamma r)^2}}\right)`$
and similarly for the N<sup>2</sup>LO correction. Note that the normalization condition, Eq. (LABEL:eq:ort1), implies a linear relation between $`c_S^{(2)}`$ and $`c_D^{(2)}`$ as well as $`c_S^{(3)}`$ and $`c_D^{(3)}`$. The total D/S ratio obtained by including the zeroth order contribution is given by
$`\eta =\eta ^{(0)}{\displaystyle \frac{1+c_D^{(2)}+c_D^{(3)}}{1+c_S^{(2)}+c_S^{(3)}}}`$ (162)
If we fix $`\eta `$ we get a relation between $`c_S`$ and $`c_D`$. The coefficients $`c_S^{(2)}`$ and $`c_D^{(2)}`$ are given by
$`c_S^{(2)}+c_S^{(3)}`$ $`=`$ $`c_1(\mathrm{})`$
$`\eta ^{(0)}\left(c_D^{(2)}+c_D^{(3)}\right)`$ $`=`$ $`c_2(\mathrm{})`$ (163)
The long distance behaviour of the integrands is well behaved since, up to inessential powers in $`r`$, one has
$`c_1^{}(r)`$ $``$ $`e^{2mr}`$
$`c_2^{}(r)`$ $``$ $`e^{2mr}`$
$`c_3^{}(r)`$ $``$ $`e^{(2m+2\gamma )r}`$
$`c_4^{}(r)`$ $``$ $`e^{(2m+2\gamma )r}`$ (164)
However, the leading short distance behaviour of the integrand is given
$`c_i^{}(r)r^{3/4}(e^{4\sqrt{2R/r}})^2\left({\displaystyle \frac{C_5}{r^5}}+{\displaystyle \frac{C_6}{r^6}}\right)r^{3/4}e^{\pm i4\sqrt{R/r}}`$ (165)
So, we expect the coefficients $`c_S`$ and $`c_D`$ to diverge if the short distance cut-off is removed, $`r_c0`$. It is unclear how this divergence might be avoided. Unlike the scattering problem in perturbation theory, where energy dependent (and hence orthogonality violating) subtractions are needed, it would be difficult to accept a bound state not normalized to unity unless one includes, besides $`pn`$ other Fock state components, such as $`pn\pi `$. This short distance analysis holds also when the $`1/r^6`$ $`\mathrm{\Delta }`$-contributions are taken into account perturbatively.
Thus, perturbation theory on the distorted OPE basis for the deuteron makes sense only as a finite cut-off theory. In the appendix A we develop further such an approach to NLO, where $`\eta `$ is an input and to NNLO, where both $`A_S`$ and $`\eta `$ should be fixed. We also show that in the cut-off theory the NLO yields tiny corrections to deuteron properties whereas the NNLO dominates. This proves that, at least perturbatively and in the absence of the $`\mathrm{\Delta }`$, the (integer) power counting to NNLO in deuteron properties is obviously not convergent. To some extent this result resembles qualitatively the findings of Ref. Fleming et al. (2000) based on the idea that OPE and TPE can be included perturbatively Kaplan et al. (1998). Of course, an amelioration of the convergence in the cut-off theory when $`\mathrm{\Delta }`$โs are included is not precluded, and deserves further investigation. However, the very need of a finite cut-off will still hold, as our analytical study shows.
In Fig. 18 we show LO, NLO and NNLO order wave functions when a finite short distance cut-off $`r_c=0.5\mathrm{fm}`$ is considered. The strong divergence of the wave function at the origin can be clearly seen.
## X Conclusions
In the present work we have extended the coordinate space renormalization of central waves in NN interaction discussed in our previous work Pavon Valderrama and Ruiz Arriola (2005a) for the OPE potential to the TPE potential. As we have stressed along the paper, the main advantage of such a framework is that the (renormalized) potential is finite everywhere except at the origin where a Van der Waals attractive singularity takes place. This suggests using a radial cut-off which provides a compact support for the short range part of the potential thus making scheme dependent contact interactions innocuous for the long range solution. Thus, there is no need to device different regularization methods for the potential and the wave functions both for bound state and scattering state solutions. As a result model independent long range correlations between NN observables can be deduced if the potential is iterated to all orders.
Important constraints can be deduced from the requirement of a small wave function in the unknown short distance region. As a consequence the boundary condition for the wave function at short distances becomes energy independent if the long range contribution to the potential is also energy independent. We stress here that such requirements, although quite natural from a physical viewpoint, may not appear obvious within the so far established EFT framework, and it would be very interesting to provide further arguments within EFT itself supporting our unconventional framework Pavon Valderrama and Ruiz Arriola (2006). Actually, we find that the singularity structure of the potential at short distances determines uniquely how many parameters must be regarded as unknown, non predictable, information. This is done in terms of short distance phases or equivalently via suitable mixed boundary conditions at the origin. Moreover, for an energy independent potential the orthogonality of wave functions precludes a possible energy dependence of the boundary conditions. In the particular cases studied in this paper, namely $`{}_{}{}^{1}S_{0}^{}`$ and $`{}_{}{}^{3}S_{1}^{}^3D_1`$ channels, we have found that besides the NNLO TPE potential parameters, one can use the S-wave scattering lengths in both channels as well as the deuteron binding energy and the asymptotic $`D/S`$ ratio of the deuteron wave functions as independent input information. The remaining scattering or bound state properties in the triplet channel are then predicted unambiguously. Based on the superposition principle of boundary conditions, we have found analytical and simple universal rational relations which clearly exhibit these features. These universal relations would be very difficult to deduce in momentum space and, moreover, are free from uncertainties attributable to finite cut-off effects. So, the cut-off has been effectively eliminated. On a numerical level, the fact that our problem is an initial value problem for the Schrรถdinger equation starting at infinity, makes possible to obtain any solution by competitive algorithms with adaptable integration steps with any prescribed accuracy. This allows to faithfully describe the short distance oscillations of the wave function. This is in contrast to the standard Lippmann-Schwinger treatments, where matrix inversion methods may eventually run into computer space limitations with a natural loss of space resolution as a side-effect. The non-trivial oscillating structure of the wave functions with ever decreasing periods of the wave functions close to the origin would actually be very difficult to reproduce within a momentum space framework.
According to our analysis, there are finite cut-off effects in previous works dealing also with TPE potentials both in coordinate as well as in momentum space. The induced corrections are larger than the experimental uncertainty of the computed observables, so that in some cases agreement with data may be clearly attributed to the choice of a finite cut-off. In our energy independent boundary condition treatment we found short distance cut-offs of about $`r_c=0.10.2\mathrm{fm}`$ to be rather innocuous. Within a Wilsonian viewpoint of renormalization, changes in the cut-off should correspond to decimation, i.e. halving, and not to linear changes in the scale. If one associates this coordinate space cut-off to a momentum space ultraviolet cut-off of $`\mathrm{\Lambda }=\pi /2r_c`$ Pavon Valderrama and Ruiz Arriola (2004d) we are dealing with an equivalent momentum scale of about $`1.53\mathrm{G}\mathrm{e}\mathrm{V}`$, much larger than the scales below $`1\mathrm{G}\mathrm{e}\mathrm{V}`$ usually employed in momentum space calculations where only linear sensitivity to changes of the cut-off was implemented. Nevertheless, it is fair to say that the calculations based on Sets III and IV provide not too large discrepancies.
As one naturally expects in a renormalized theory, errors are dominated by uncertainties in the input data, and not by cut-off uncertainties. Indeed, we seem to reach a limit in the accuracy of the predictions, paralleling the findings in ChPT for mesons at the two loop level. At the OPE level, one can predict bound state and scattering properties in the singlet $`{}_{}{}^{1}S_{0}^{}`$ and triplet $`{}_{}{}^{3}S_{1}^{}^3D_1`$ channels solely from the deuteron energy and the $`{}_{}{}^{1}S_{0}^{}`$ scattering length. At the TPE level, one needs not only the additional chiral constants $`c_1`$, $`c_3`$ and $`c_4`$ but also the triplet S-wave scattering length and the asymptotic $`D/S`$ deuteron ratio. Although the TPE central values predictions improve, the induced TPE errors turn out to be larger than the OPE uncertainties. In fact, this large uncertainties make that, within errors, the TPE calculation becomes compatible with experimental data at the $`1\sigma `$ level. This suggests that in order to see in a statistically significant sense other effects, such as electromagnetic, relativistic and three pion effects one must first improve on the input data. Otherwise, predictive power is lost. Nevertheless, given the finite cut-off effects detected in previous works, the role of these corrections beyond TPE should be reanalyzed within the present approach.
One of the important consequences of our treatment is that the chiral constants $`c_1`$, $`c_3`$ and $`c_4`$ can be determined from low energy data and deuteron properties. Specifically, we have used the singlet and triplet effective ranges as well as the asymptotic S-wave deuteron wave function to $`c_1`$, $`c_3`$ and $`c_4`$ with errors varying all input data within their experimental uncertainties. The decision on what set of data should be used to pin down the chiral coefficients is not entirely trivial, because it should become clear which hypothesis we want to verify or to refute. The absence of cut-off effects makes this test cleaner; we just check whether the TPE potential holds from zero to infinity. Obviously, this cannot be literally true, but one expects that at low energies other short range effects can be considered negligible. Let us remind that error analysis within NN calculations was only carried out in a large scale partial wave analysis to data in Ref. Rentmeester et al. (1999). The determinations of chiral constants based on a fit to NN databases Stoks et al. (1993, 1994); Arndt et al. (1994) for phase shifts lack any error estimates because the databases themselves are treated as errorless. The determination of chiral constants from peripheral waves has similar drawbacks. From the chiral theory point of view we see that it is possible to determine these parameters precisely in the regime where we trust the theory most, namely in the description of low energy NN data. A fit becomes possible, and the values it yields only differ by $`2\sigma `$ with the determination from $`\pi N`$ data. We do not exclude that our values for the chiral constants may eventually spoil the successful overall fit of phase shifts in all channels presented in the past, after all renormalization has been carried out. If so, the situation on the effectiveness of effective field theory would be in a less optimistic shape than assumed hitherto. A preliminary analysis of the problem shows what Van der Waals coefficients in the TPE potential correspond to attractive short range interactions, and hence what phase shifts are completely determined in terms of coupled channel scattering lengths. This issue is very relevant and would require a detailed channel by channel analysis and renormalization, taking as input the scattering lengths documented in our previous work Pavon Valderrama and Ruiz Arriola (2004c) and integrating in from large distances along the lines of the present approach. Full details are reported elsewhere Pavon Valderrama and Ruiz Arriola (2005b).
Nevertheless, despite of the good convergence in the $`{}_{}{}^{1}S_{0}^{}`$ channel for LO, NLO and NNLO calculations, we have noted a difficulty for the triplet $`{}_{}{}^{3}S_{1}^{}^3D_1`$ channel at NLO of the potential. In contrast to dimensional power counting expectations one cannot use the scattering length $`\alpha _0`$, the effective range $`r_0`$ and $`\alpha _{02}`$ as arbitrary input parameters at NLO in the potential (one could equally take $`\gamma `$, $`\eta `$ and $`\alpha _0`$) but they are entirely predicted from the potential as required by finiteness of the phase-shifts. Otherwise, the scattering amplitude diverges, as we have shown. We have also seen that even if one assumes a finite value of the cut-off the NLO is worse than the LO, suggesting that the problem may indeed be related to the power counting on the long distance potential. Remarkably, these parameters must be fixed at NNLO where, according to the standard approach, no further low energy parameters should be fixed. This mismatch in orders can be understood if one considers the $`N\mathrm{\Delta }`$ splitting to be a small parameter, making much of the NNLO contributions to the potential become NLO ones, because $`c_3`$ and $`c_4`$ would be order minus one. In such a case, our interpretation goes hand in hand with the standard approach; one needs three independent low energy parameters at NLO in this counting. The consequences of this $`\mathrm{\Delta }`$-counting to higher orders within the context of renormalization will be explored elsewhere. Of course, we should point out that despite the rather tantalizing description achieved at NNLO the existence of a consistent power counting guaranteeing the success of the present approach to all orders remains to be proved. A key ingredient of such a power counting would be the correct incorporation of all long range physics. Apparently, within Weinbergโs power counting the NLO in the deuteron channel misses important contributions.
Finally, we have analyzed the consequences of a perturbative expansion of TPE effects taking the OPE results as a zeroth order approximation as suggested recently Nogga et al. (2005). Our non-perturbative calculations based on iterating a perturbative potential to all orders exhibit unequivocal non-analytic dependence on the expansion parameter, due to the singular character of the chiral potentials at the origin. This is equivalent to a non-integer enhancement of the power counting $`\lambda ^\alpha `$ with $`0<\alpha <1`$ in the potential, and it would be interesting to know the general rules of such a counting a priori Calle Cordรณn et al. (2006). Thus, perturbation theory based on standard power counting becomes divergent and can only yield finite results at the expense of introducing more perturbative counterterms than are needed in a non-perturbative treatment. This is just a manifestation of the fact that singular potentials require infinite counterterms in perturbation theory, while only a few ones are needed non-perturbatively. Specifically, our analysis shows that it would be necessary to include at least 4 counterterms for the singlet $`{}_{}{}^{1}S_{0}^{}`$ and 6 counterterms for the triplet $`{}_{}{}^{3}S_{1}^{}^3D_1`$ channel at NNLO. This proliferation of counterterms is expected to occur also in other partial waves because the singularity of the potential dominates over the centrifugal barrier at short distances. In the $`{}_{}{}^{1}S_{0}^{}`$ channel, we have seen that adding more counterterms in fact worsens the results for the effective range expansion parameters. In contrast, our non-perturbative calculations are based on just 1 and 3 counterterms respectively. The good quality of our results suggests that our choice of less counterterms cannot be refuted on the basis of phenomenology. In the deuteron case we have made a calculation to NNLO in perturbation theory. Our analysis shows that such a perturbative approach only makes sense if a finite cut-off is introduced. In any case, the cut-off theory has less predictive power, does not provide a better phenomenological description of the deuteron than our non-perturbative renormalized results and is non-convergent since NNLO corrections are numerically much larger (two or three orders of magnitude) than NLO ones, despite being parametrically small. In our view this is a perturbative manifestation of the short distance dominance which has been unveiled non-perturbatively. In addition, the difficulties faced by a perturbative treatment are simply absent in the non-perturbative approach.
One of the main goals of nuclear physics is the determination of the nucleon-nucleon interaction. From a theoretical viewpoint the disentanglement of such an interaction in terms of pion exchanges based on chiral symmetry requires dealing with non-trivial and, to some extent, unconventional non-perturbative renormalization issues in the continuum, but it is crucial because it shows our quantitative understanding of the underlying theory of quarks and gluons in the chirally symmetric broken phase. Our results also show that the singular chiral Van der Waals forces are not necessarily spurious and inconvenient features of the chiral potential. Instead, as we have shown, the singularities alone in conjunction with renormalization ideas explain much of the observed $`S`$waves phase shifts with natural values of the chiral constants, and provide an appealing physical picture. In this regard, it is interesting to realize that based on the analogy with molecular systems, which also exhibit a long range Van der Waals force, the liquid drop model was formulated more than 60 years ago. Chiral dynamics may provide not only a closer analogy and perhaps more quantitative insights into the hydrodynamical and thermodynamical properties of nuclei but also a theoretical justification from the underlying theory of strong interactions.
###### Acknowledgements.
One of us (E.R.A.) thanks M. Rentmeester, R. Machleidt, E. Epelbaum, N. Kaiser and G. Colangelo for useful correspondence. We thank them and also R. Higa and A. Nogga for discussions and D. Phillips for stressing the role of the $`\mathrm{\Delta }`$. We thank J. Nieves for reading an early version of the manuscript. This work is supported in part by funds provided by the Spanish DGI with grant no. FIS2005-00810, Junta de Andalucรญa grant no. FM-225 and EURIDICE grant number HPRN-CT-2003-00311.
## Appendix A The deuteron in OPE-distorted perturbation theory with a cut-off to NNLO
In this appendix we illustrate the situation discussed in Sect. IX by solving numerically the set of perturbative equations (LABEL:eq:pert\_first). As we have mentioned such a calculation makes only sense within a finite cut-off scheme.
In practice, we integrate from large distances ($`25\mathrm{f}\mathrm{m}`$) with the conditions specified by Eq. (LABEL:eq:long\_pert) with some prescribed values of $`c_S`$ and $`c_D`$ <sup>20</sup><sup>20</sup>20This is numerically more efficient and stable procedure than a direct use of the explicit expressions Eq. (163) involving determinants.. This can be advantageously done using the superposition principle of boundary conditions, Eq. (60), yielding in perturbation theory
$`u_\gamma (r)`$ $`=`$ $`u_\gamma ^{(0)}(r)+u_\gamma ^{(2)}(r)+u_\gamma ^{(3)}(r)+\mathrm{}`$
$`w_\gamma (r)`$ $`=`$ $`u_\gamma ^{(0)}(r)+u_\gamma ^{(2)}(r)+u_\gamma ^{(3)}(r)+\mathrm{}`$ (166)
At LO the wave function can be written as
$`u_\gamma ^{(0)}(r)`$ $`=`$ $`u_S^{(0)}(r)+\eta ^{(0)}u_D^{(0)}(r)`$
$`w_\gamma ^{(0)}(r)`$ $`=`$ $`w_S^{(0)}(r)+\eta ^{(0)}w_D^{(0)}(r)`$ (167)
and $`\eta ^{(0)}`$ is determined from the regularity condition at the origin Pavon Valderrama and Ruiz Arriola (2005a). At LO the normalization factor is
$`{\displaystyle \frac{1}{(A_S^{(0)})^2}}={\displaystyle _0^{\mathrm{}}}๐r(u_\gamma ^{(0)}(r)^2+w_\gamma ^{(0)}(r)^2)`$ (168)
The NLO and NNLO contributions are
$`u_\gamma ^{(2)}(r)`$ $`=`$ $`c_S^{(2)}u_S^{(2)}(r)+\eta ^{(0)}c_D^{(2)}u_D^{(2)}(r)`$
$`w_\gamma ^{(2)}(r)`$ $`=`$ $`c_S^{(2)}w_S^{(2)}(r)+\eta ^{(0)}c_D^{(2)}w_D^{(2)}(r)`$
$`u_\gamma ^{(3)}(r)`$ $`=`$ $`c_S^{(3)}u_S^{(3)}(r)+\eta ^{(0)}c_D^{(3)}u_D^{(3)}(r)`$
$`w_\gamma ^{(3)}(r)`$ $`=`$ $`c_S^{(3)}w_S^{(3)}(r)+\eta ^{(0)}c_D^{(3)}w_D^{(3)}(r)`$ (169)
The advantage is that the functions appearing here only depend on the potential and the deuteron binding energy, whereas the coefficients must be determined by some additional conditions. In the first place, normalization to NLO and NNLO requires orthogonality of the wave functions to the LO solution,
$`0`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}๐r(u^{(0)}(r)u^{(2)}(r)+w^{(0)}(r)w^{(2)}(r))`$
$`0`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}๐r(u^{(0)}(r)u^{(3)}(r)+w^{(0)}(r)w^{(3)}(r))`$ (170)
This implies the couple of linear relations <sup>21</sup><sup>21</sup>21For instance, at $`r_c=0.5\mathrm{fm}`$ we get $`c_D^{(2)}=5.865c_S^{(2)}`$ and $`c_D^{(3)}=5.545c_S^{(3)}`$
$`\eta ^{(0)}{\displaystyle \frac{c_D^{(2)}}{c_S^{(2)}}}`$ $`=`$ $`{\displaystyle \frac{_0^{\mathrm{}}๐r(u^{(0)}(r)u_S^{(2)}(r)+w^{(0)}(r)w_S^{(2)}(r))}{_0^{\mathrm{}}๐r(u^{(0)}(r)u_D^{(2)}(r)+w^{(0)}(r)w_D^{(2)}(r))}}`$
$`\eta ^{(0)}{\displaystyle \frac{c_D^{(3)}}{c_S^{(3)}}}`$ $`=`$ $`{\displaystyle \frac{_0^{\mathrm{}}๐r(u^{(0)}(r)u_S^{(3)}(r)+w^{(0)}(r)w_S^{(3)}(r))}{_0^{\mathrm{}}๐r(u^{(0)}(r)u_D^{(3)}(r)+w^{(0)}(r)w_D^{(3)}(r))}}`$
Further relations can be obtained by imposing renormalization conditions. Note that the required number of conditions increases with the order. This is similar in spirit to the procedure of adding more counterterms for the scattering problem discussed in Sect. IX. For instance, using the perturbative expansion for $`A_S`$ and $`A_D`$
$`A_S`$ $`=`$ $`A_S^{(0)}\left(1+c_S^{(2)}+c_S^{(3)}+\mathrm{}\right)`$
$`A_D`$ $`=`$ $`A_D^{(0)}\eta ^{(0)}\left(1+c_D^{(2)}+c_D^{(3)}+\mathrm{}\right)=\eta A_S`$
In practice, we use a short distance cut-off $`r_c`$ for the NLO and NNLO contributions only. Deuteron properties can be written to NNLO as follows
$`r_m`$ $`=`$ $`r_m^{(0)}+c_S^{(2)}r_m^{(2,S)}+\eta ^{(0)}c_D^{(2)}r_m^{(2,D)}`$ (173)
$`+`$ $`c_S^{(3)}r_m^{(3,S)}+\eta ^{(0)}c_D^{(3)}r_m^{(3,D)}+\mathrm{}`$
$`Q_d`$ $`=`$ $`Q_d^{(0)}+c_S^{(2)}Q_d^{(2,S)}+\eta ^{(0)}c_D^{(2)}Q_d^{(2,D)}`$ (174)
$`+`$ $`c_S^{(3)}Q_d^{(3,S)}+\eta ^{(0)}c_D^{(3)}Q_d^{(3,D)}+\mathrm{}`$
where the potential contributions have explicitly been factored out. The numerical solution requires some care, due to the short distance instabilities and oscillations. This requires using an adaptive grid to optimize the convergence. Since solutions of different orders must be mixed in the evaluation of the orthogonality conditions, Eq. (LABEL:eq:cS\_cD\_ort), and observables, Eq. (174), we solve all LO, NLO and NNLO equations simultaneously to provide all functions on the same grid.
At NLO and fixing $`r=r_c`$ we demand the experimental value of $`\eta `$, from Eq. (162) and Eq. (LABEL:eq:cS\_cD\_ort). This way, a solution which we denote by $`(c_S^{(2)}|_{\mathrm{NLO}},c_D^{(2)}|_{\mathrm{NLO}})`$ can be obtained. From there we can obtain deuteron properties to NLO, as a function of the perturbative cut-off $`r_c`$ In Fig. 19 we show the dependence of $`A_S`$, $`r_m`$, $`Q_d`$ and $`p_d`$ on $`r_c`$. As we see the NLO correction is tiny and stable for $`r_c>0.2\mathrm{fm}`$. At NNLO we fix $`\eta `$ and $`A_S`$. The solution is now $`(c_S^{(2)}|_{\mathrm{NNLO}},c_D^{(2)}|_{\mathrm{NNLO}})`$ and $`(c_S^{(3)}|_{\mathrm{NNLO}},c_D^{(3)}|_{\mathrm{NNLO}})`$. Note that in general the NLO coefficients $`c_S^{(2)}`$ and $`c_D^{(2)}`$ must be readjusted. In this case the correction is much larger than the NLO case (see Fig. 19) and the cut-off dependence is stronger due to the the $`1/r^6`$ singularity of the NNLO potential. As a curiosity, we mention that at short distances worry-some negative D-wave probabilities show up below $`r_c=0.17\mathrm{fm}`$ at NNLO, a spurious feature which can only take place in perturbation theory and sets a unitarity bound on the short distance perturbative cut-off. Numerical results are provided in Table 8 for Set IV. Typically, we find that results do not depend dramatically on the chosen chiral couplings . We take $`r_c=0.5\mathrm{fm}`$ as a standard choice. As we see, finite cut-off perturbation theory does not work better than our non-perturbative results of Sect. V, and in fact requires one more counterterm. Actually, this is a perturbative indication that NNLO is more important than NLO, casting doubts on the convergence of the approach.
Finally, we have checked that taking the LO to be the full OPE potential, $`๐_{1\pi }`$ and the perturbation to be $`๐_{2\pi }^{(2)}+๐_{2\pi }^{(3)}`$ as a whole and keeping the cut-off $`r_c>0.1\mathrm{fm}`$ does not change the results significantly. Actually, the perturbative result does not account for the value obtained non-perturbatively, despite the apparent linear behaviour observed when changing numerically the scaling parameter $`\lambda `$ in the region $`\lambda 0.1`$ (see Fig. 17). This supports our conclusion that perturbation theory does not compute the slope of $`A_S(\lambda )`$ at the origin. In addition, even if we disregard the divergence by introducing a cut-off, the perturbative calculation does not account for the non-perturbative renormalized result.
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# The Supernova Rate-Velocity dispersion Relation in the Interstellar Medium
## 1 INTRODUCTION
It is now a well established fact that the interstellar medium (ISM) in galactic disks is turbulent (Larson 1981; Scalo 1987; Dickey & Lockman 1990; Elmegreen & Scalo 2004). In most spiral galaxies, and after correcting for instrumental effects, the vertical velocity dispersion derived essentially from H i observations is observed to vary radially from $`1215`$ km s<sup>-1</sup> in the central parts to $`46`$ km s<sup>-1</sup> in the outer parts (van de Kruit & Shostak 1982 for NGC 3938; Shostak & van der Kruit 1984 for NGC 628; Dickey, Hanson & Helou 1990 for NGC 1058; Kamphuis & Sancisi 1993 for NGC 6964; Rownd, Dickey & Helou 1994, for NGC 5474; Meurer, Mackie & Carignan 1994 and Meurer et al. 1996 for NGC 2915; de Blok & Walter 2005 for NGC 6822). In most cases, the latter values exceed those expected by the thermal broadening of the H i emission line. As an example, the radial dependence of the velocity dispersion in NGC 1058 is shown in Fig. 1. Fig. 2 shows the behavior of the characteristic velocity dispersion of a sample of galaxies when plotted versus their characteristic star formation rate. The galaxies shown in Fig. 2 correspond to a sample where both the velocity dispersion and the star formation rate were available in the literature. Similarly to the observed velocity dispersion radial profiles, the data plotted in Fig. 2 argues for the existence of a minimum level of turbulence in galactic disks. The case of NGC 2915 is particularly intriguing as it has an extended H i disk with a constant velocity dispersion, even in the outer parts, where stars do not form (Meurer et al. 1996).
Many physical processes might contribute to the driving of the turbulence in the ISM, acting on different scales and injecting different amounts of kinetic energy into the medium. These driving mechanism can be related to stellar activity (e.g., ionizing radiation from H ii regions, jets from young stellar objects, stellar winds, supernova explosions), or to gas hydrodynamical or magnetohydrodynamical instabilities (e.g., thermal and gravitational instabilities, magnetorotational instability, Kelvin-Helmotz, Rayleigh-Taylor and Parker instabilities) or, eventually, to the complex interaction between these processes. Mac Low & Klessen (2004) showed, using simple analytical estimates, that the global energy input into the ISM of the Galaxy from supernova explosions (SNe) is much larger than from any of the above mentioned processes. However, their estimates are global ones which do not compare the efficiency of each process at different galactic radii and the possible interactions between two or several processes. Other, yet poorly explored, sources of kinetic energy injection into the ISM of some galaxies might reside in perturbations occurring on galactic scales (e.g., frequent minor mergers, tidal interactions with satellite galaxies and ram pressure effects).
In the past decade, some of the above mentioned energy injection mechanisms into the ISM have been investigated by means of numerical simulations. Wada & Norman (1999,2001), Wada, Spaans & Kim (2000) and Wada, Meurer & Norman (2002) investigated the evolution of thermal and gravitational instabilities (TI and GI, respectively) in galactic disks. Wada, Meurer & Norman (2002) showed, using high resolution two-dimensional simulations of an NGC 2915 like disk that the turbulent energy spectra can be maintained in a quasi-stationary state in the absence of stellar feedback. In the latter simulations, the kinetic energy decay and radiative cooling are compensated by the interaction of the galactic shear with the gas self-gravity. Another indication that energy might be injected into the ISM of low star forming galaxies on large scales before cascading to smaller scales via thermal and gravitational instabilities is provided by the analysis of their H i gas morphology. By comparing the H i emission map of Holmberg II (Ho II) to synthetic H i maps of driven turbulence which include also cooling, heating and gravity, Dib & Burkert (2005) found that kinetic energy is injected into the ISM of Ho II on a scale of $`6`$ kpc, which is much larger than the scale implied by SN driving. Galactic shear might play a role in supporting turbulence in Ho II, however, the comet shape structure of the H i gas in Ho II (Bureau et al. 2004) suggests that turbulence might be also partially induced by the effect of tidal interactions with Ho IIโs two satellite galaxies.
The coupling of the galactic shear to magnetic fields can trigger the magnetorotational instability (MRI) (Balbus & Hawley 1991). MRI has been the scenario invoked by Selwood & Balbus (1999) to explain the constancy of the velocity dispersion in the outer parts of NGC 1058. The main argument of Selwood & Balbus (1999) against a SN-driven ISM in the outer regions of NGC 1058, in addition to the fact that column densities are low and star formation is inefficient, is the observed uniformity of the H i velocity dispersions (Dickey, Hanson & Helou 1990) whereas H ii regions in NGC 1058 are primarily observed in narrow spiral arms (Ferguson et al. 1998). Local, three-dimensional simulations with a single phase medium by Kim, Ostriker & Stone (2003) of the MRI result in velocity dispersions of the order $`1.63.2`$ km s<sup>-1</sup>. Global one-phase medium simulations by Dziourkevitch et al. (2004) and Dziourkevitch (2005) show that some amount of turbulent motions can be created by the MRI in galactic disks. However, this turbulence is mostly located in the inner parts of the disk (the inner 4 kpc) and the velocity dispersion does not exceed in that case $`3.54`$ km s<sup>-1</sup>, dropping quickly to very small values at larger radii. Two and three-dimensional MRI simulations in a medium affected by the TI but with no stellar feedback have been presented by Piontek & Ostriker (2004,2005). In the three-dimensional models, the latter authors show that the velocity dispersion depends on the average density of the medium (larger dispersions at lower densities). At the lowest density they have considered (0.25 cm<sup>-3</sup>), the one-dimensional velocity dispersion is found to be 2.2, 2.2 and 1.1 km s<sup>-1</sup> for the three directions of the box, respectively ($``$ 4.5 km s<sup>-1</sup> for the three components if thermal broadening is not subtracted). The latter values are a factor of 3-6 lower than the observational values ($``$ 1.5-2 if thermal broadening is not subtracted). Recent simulations of the MRI in a multiphase medium with star formation included show that the MRI might be completely suppressed/overwhelmed by stellar feedback (M. Korpi, private communication).
Another major source of energy input into the ISM is the stellar energy feedback, particularly from massive stars. The latter can be delivered to the ISM in the form of ionizing radiation and stellar winds from O and B stars (Kessel-Deynet & Burkert 2003) and from clustered or field supernova explosions (McKee & Ostriker 1977). Two and three-dimensional numerical models of SN explosions models in the ISM have been presented in the literature (Rosen & Bregman 1995; Korpi et al. 1999a,1999b; Gazol-Patiรฑo & Passot 1999; de Avillez 2000; de Avillez & Berry 2001; Avila-Reese & Vรกzquez-Semadeni 2001; Kim et al. 2001; Kim 2004; de Avillez & Breitschwerdt 2004,2005; Slyz et al. 2005; Mac Low et al. 2005). In these models, the authors have focused on problems like the evolution of SNe bubbles and their outburst through the galactic disk, the halo-disc interaction, the vertical scale heights and volume filling factors of the different gaseous phases and the effects of supernova explosions on the Galactic dynamo (Ferriรจre 1992a,1992b,1998a,1998b,2000). However, the existence of a correlation between the SN rate and the velocity dispersion of the gas has not been investigated so far. A different interpretation of the velocity dispersion in the outer regions of galactic disks is discussed by Schaye (2004), which argues that if the surface density in the outer parts of the disk falls below a critical value (i.e., column density $`Log(N_H)<27.5`$ cm<sup>-2</sup>), the medium can be efficiently heated to a temperature of $`T8000`$ K by a background of ultraviolet (UV) ionizing radiation, both of galactic and extragalactic origin. Thus, thermal broadening could account for the observed level of line broadening in the outer parts of galactic disks, leaving little room for a dynamical, turbulent component. This hypothesis can be appropriately tested if the level of UV radiation, particularly the extragalactic component, and the temperature radial profiles in galaxies would be better constrained from observations.
The aim of this paper is to assess how much of the velocity dispersion observed in the ISM of galaxies can be due to SN feedback for various values of the SN rate and feedback efficiency. A particular point of interest is to understand the constancy of the velocity dispersion in galactic disks at different radii where the star formation rate is expected to decrease with increasing radius from the galactic center as predicted by many empirical star formation laws (Schmidt 1959,1963; Kennicutt 1998a,1998b; Dopita & Ryder 1994; Prantzos & Silk 1998; see also Li et al. 2005a,2005b who finds, as an explanation of the global Schmidt laws observed in galaxies, that the star formation rates correlate with the efficiency of gravitational instability). This paper is organized as follows. In ยง 2, we describe our models and the relevant parameters. In ยง 3, we describe how synthetic observations are derived. The velocity dispersion dependence on the feedback efficiency, supernova rate and average gas number density is presented and discussed in ยง 4, ยง 5 and ยง 6 respectively. Detailed comparison to the observations is performed in ยง 7. In ยง 8 the need for improved numerical models is critically reviewed and in ยง 9, we summarize our results and conclude.
## 2 THE MODEL
In order to understand how the SN rate and energy feedback efficiency affect the velocity dispersion of the gas, we resort to a simple numerical model in which the vertical stratification, galactic rotation, magnetic fields and the gas self-gravity are not included. We use the ZEUS-3D code (Stone & Norman 1992a,b) to solve the equations (mass, momentum and internal energy conservation) of ideal gas dynamics. We simulate a 1 kpc<sup>3</sup> volume of the ISM with a grid resolution of 128<sup>3</sup>. Periodic boundary conditions are imposed in the three directions. In most simulations and if not specified otherwise, the initial density field is homogeneous with a number density of $`\overline{n}=0.5`$ cm<sup>-3</sup>. The velocity and temperature are everywhere zero and 10<sup>4</sup> K, respectively. Radiative cooling of the gas is included by directly interpolating in the the solar metallicity cooling curves of Dalgarno & McRay (1972) in the temperature range of \[100 K,10<sup>4</sup> K\] and of the more recent data of Sutherland & Dopita (1993) for the temperature range of \[10<sup>4</sup> K-10<sup>8.5</sup> K\]. The gas is not allowed to cool below the minimum temperature $`T_{min}=100`$ K. The maximum temperatures at the explosion sites reach values of $`6070\times 10^6`$ K which makes the hot gas fall in the hot stable regime. A polytropic equation of state with a specific heat ratio of $`5/3`$ is used. This proves to be justified because of the lower cutoff temperature of $`100`$ K that limits over-densities to a few tens of cm<sup>-3</sup> and therefore the gas remains mostly monoatomic.
Stellar feedback is modeled as resulting from type II SN explosions only and the energy is injected into the ISM instantaneously. The total energy of each explosion is taken to be $`E_{SN}=10^{51}`$ erg (Chevalier 1977; Abbott 1982; Woosley & Weaver 1986 and Heiles 1987). Only a fraction of the total energy is transferred to the ISM in the form of thermal energy. This defines the feedback efficiency parameter $`ฯต`$. The energy is injected into $`3^3`$ cells around the central explosion cell following a Gaussian profile. The energy each neighboring cell receives is weighted by the square of its distance to the central cell in order to insure a better isotropy of the explosion. The site where a new SN explosion occurs is chosen randomly under the condition that the local density $`n`$ is such that $`n\overline{n}`$. This assumption leads to a more realistic fraction of clustered SN explosions. The time interval between two consecutive SN explosions is given by $`\mathrm{\Delta }t_{SN}=1/\eta `$, where $`\eta `$ is the SN explosion rate. However, when the time interval between two consecutive SN explosions become shorter than the CFL (Courant-Friedrisch-Levy) time step, more than one SNe are detonated simultaneously at different locations of the grid. The number of SNe detonated in that case is taken to be the closest integer to the ratio $`dt/\mathrm{\Delta }t_{SN}`$, where $`dt`$ is the CFL time step. We define a Galactic SN rate ($`\eta _G`$) of 2.58 $`\times 10^4`$ yr<sup>-1</sup> kpc<sup>-3</sup> assuming a Galactic radius of 15.5 kpc and a scale height for type II SN of 90 pc (Miller & Scalo 1979). We use a frequency of 1/57 yr<sup>-1</sup> (Capellaro, Evans & Turatto 1999) which is smaller than earlier estimates of 1/50 yr<sup>-1</sup> by van den Bergh & McClure (1990) and 1/37 yr<sup>-1</sup> by Tammann, Lรถffler & Schrรถder (1994). Note that in our simulations, heating by a UV background radiation is not accounted for. It has not been included because no clear recipe exists that correlates the amount of UV heating that should be added with the different SN rates. One possibility is to distribute a certain fraction of the energy associated with a given SN rate on the total grid. However, at low SN rates, associated with low density environments in the outer parts of galactic disks, the disk might be more easily heated by incident UV photons of extragalactic origin which intensity is not well determined (see discussion in Schaye 2004).
In order to test the effects of numerical resolution, we explode a single supernova in a medium of initial temperature of $`10^4`$ K, initial homogeneous and uniform density of $`0.5`$ cm<sup>-3</sup> on a grid representing a physical scale of 1 kpc. The test is similar to the one presented by Mac Low et al. (2005). However the latter authors, aside from having a smaller simulations box (200 pc), assumed the SN remnant to evolve in a medium of negligible pressure (i.e., temperature of 10 K and average density of 0.1 cm<sup>-3</sup>). Fig. 3 displays the time evolution of the shell position R<sub>sh</sub>, measured as being the position of the density peak, in simulations with numerical resolutions of $`64^3`$, the fiducial resolution of $`128^3`$, and $`256^3`$. Clearly, a resolution of $`64^3`$ is not sufficient to resolve properly the dynamics of a single SN remnant. On the other hand, particularly after the first $`0.1`$ Myrs in the lifetime of the SN remnant, the discrepancy in the position of the shell at the resolutions of $`128^3`$ and $`256^3`$ is of the order of $`10\%`$ and does not exceed the value of $`20\%`$. We conclude that our fiducial resolution of $`128^3`$ is good enough for the purpose of studying the global dynamical effects of the energetic input of SN explosions into the ISM.
## 3 ANALYSIS AND DERIVATION OF THE OBSERVABLES
The simulations are evolved until kinetic energy reaches a stationary value. Fig. 4 shows the evolution of kinetic energy in a number of simulations. The equilibrium value for the kinetic energy is reached when the dissipation equals the amount of injected kinetic energy. The medium acquires kinetic energy from the SNe explosion-induced thermal pressure gradients, and the thermal pressure gradients associated with TI which occurs in the dense expanding shells. In most simulations, the equilibrium of thermal energy is also reached except for the simulations with a high supernova rate ($`2.5\eta _G`$). In the latter simulations, the overlap radii between supernova remnants is very small, and all the gas continues to heat up as more and more energy is injected into the system. Fig. 5 and Fig. 6 show snapshot two-dimensional cuts for models with $`(\eta /\eta _G,ฯต)=(0.1,0.25)`$ and $`(1,0.25)`$, respectively. In the simulations with the lower rates, larger and denser clouds are able to form under the effect of TI, before being eventually dispersed by a next generation local SN explosion.
We calculate the velocity dispersion of the gas in two complementary ways. In the first method, we evaluate the characteristic mass weighted velocity $`v_c`$ from the three-dimensional data as
$$v_c=\sqrt{\frac{\mathrm{\Sigma }_{i=1}^{n_{cells}}m_i|v_i|^2}{\mathrm{\Sigma }_{i=1}^{n_{cells}}m_i}},$$
(1)
where the index $`i`$ runs over the number of cells in the simulation box. An average is made over the last 5 Myrs (5 values) in each simulations in order to smooth for time fluctuations which is particularly useful for the low SN rate simulations. On the other hand, following a more observational approach, the one-dimensional velocity dispersion $`\sigma `$ is obtained by fitting the mass-weighted line of sight velocity profile. Intensity is assumed to be proportional to the mass along the line of sight. This is particularly true in the case of the H i line (Rholfs & Wilson 1996). The velocity profile is then normalized to its maximum value. A second velocity dispersion which we call $`\sigma _{\mathrm{H}\mathrm{i}}`$ is obtained by fitting a velocity profile where only cells which have temperatures $`12000`$ K and number densities $`n0.25`$ cm<sup>-3</sup> have been accounted for, thus mimicking the velocity profile of an H i emission line. We have tested the dependence of the fit parameters on the size of the adopted velocity bin. The relevant parameter (i.e., width of the velocity profile) is practically unchanged as we vary the velocity bin size from 0.1 km s<sup>-1</sup> up to 2.5 km s<sup>-1</sup>, only the fit-error on the parameters changes, remaining however very close to the value of the spectral resolution. The results we will show correspond to a spectral resolution of 1 km s<sup>-1</sup>. The latter value is characteristic of single-dish radio telescopes (e.g., Effelsberg radio telescope, Green Bank radio telescope) and nears the spectral resolutions obtained with the VLA (Very Large Array) $`2.5`$ km s<sup>-1</sup> which will be enhanced when the EVLA (Extended Very Large Array) becomes operational. Here also, for each simulation, an averaging over the last 5 Myrs has been performed (5 estimates). Errors on $`\sigma `$ and $`\sigma _{\mathrm{H}\mathrm{i}}`$ are average values of the individual errors derived from the parameters of the fit functions whereas the error on $`v_c`$ is simply a statistical error over the 5 estimates.
We have attempted to fit the line of sight velocity profiles with several functional distributions, namely, a Gaussian, a Lorentzian, a Moffat profile (modified Gaussian) and a Voigt profile (see Lang 1980 for the mathematical definition of each function). In all cases, the Gaussian, Lorentzian and Moffat fit functions yield the same dispersion value which has a slight different meaning in each case (see Lang 1980 for details). Fitting with the Voigt profile works only in a limited number of cases, but in those case proves to be a better fit of the profiles. Examples are shown in Figs. 7-10. Fig. 7 and Fig. 8 show the 1 km s<sup>-1</sup> binned total gas and H i gas line profiles for the model with $`(\eta /\eta _G,ฯต)=(1,0.25)`$, respectively, whereas Fig. 9 and Fig. 10 show the same profiles for the model with $`(\eta /\eta _G,ฯต)=(0.1,0.25)`$. Though Gaussian fitting proves to be quite satisfactory, it is worth mentioning at this stage that our simulated velocity profile have wings that are slightly broader than those associated to Gaussian functions. This non-Gaussianity of the line profiles has been already pointed out by Ricotti & Ferrara (2002) in their Monte-Carlo models of the ISM dynamics. In a realistic and turbulent ISM, one expects that the turbulent contribution to the line broadening would be described with a Lorentzian and the thermal broadening of the line by a Gaussian. The convolution of the Gaussian and Lorentzian function results in a Voigt profile. However, even when we introduce a thermal broadening component to the line profile (in ยง 7), the Voigt profile fails to systematically fit the line profiles of the different models.
Since energy injection in our simulations is discrete, a similar value $`\eta \times ฯต`$ with different permutations of $`\eta `$ and $`ฯต`$ does not yield similar results because of the differences in the interactions of the expanding shells and non-linear development of TI. Ideally, a full investigation of the two-dimensional parameter space would be necessary, however, this would lead us beyond our current computational capabilities. At this stage, we have taken an intermediate approach and have explored the effects of $`\eta `$ and $`ฯต`$ independently by fixing one parameter and varying the other.
## 4 THE EFFECT OF THE FEEDBACK EFFICIENCY
We performed a first set of simulations for which we varied the supernova feedback efficiency between 0.05 and 1 for a constant value of the supernova rate, $`\eta =0.1\eta _G`$. Fig. 11 shows the dependence of the characteristic velocity $`v_c`$, line of sight velocity dispersion $`\sigma `$ and H i line of sight velocity dispersion $`\sigma _{\mathrm{H}\mathrm{i}}`$ on $`ฯต`$. Assuming that all of the three values would be zero if the feedback efficiency is zero, this would mean that $`v_c`$, $`\sigma `$ and $`\sigma _{\mathrm{H}\mathrm{i}}`$ would rapidly increase with increasing $`ฯต`$, followed thereafter by a slower increase at larger values of $`ฯต`$. As a velocity has the dimensions of the square root of an energy, we have attempted to fit $`v_c`$, $`\sigma `$ and $`\sigma _{\mathrm{H}\mathrm{i}}`$ by functions of the form $`A\sqrt{ฯต}`$. Fits with these functional forms are over-plotted to the data in Fig. 11. The values of the fit parameter $`A`$ are $`17.15\pm 0.59`$ km s<sup>-1</sup>, $`7.91\pm 0.36`$ km s<sup>-1</sup> and $`6.49\pm 0.44`$ km s<sup>-1</sup> for $`v_c`$, $`\sigma `$ and $`\sigma _{\mathrm{H}\mathrm{i}}`$, respectively. In all three curves the simulations are noticeably higher than the estimate of the fit functions for the small values of $`ฯต`$. This is plausibly an indication that a second mechanism is responsible for generating kinetic energy in the medium on top of the direct supernova driving in the regime of low energy injection. TI is likely to play an important role in that regime as clouds have time to condense further, thus enhancing the effects of TI, before being destroyed by a local, next generation SN explosion. We will particularly focus on the role of TI in the low energy injection regime in the next section in which we investigate the dependence of the velocity dispersion on the SN rate and for which we have more available simulations.
The total feedback energy (kinetic+thermal) deposited in the ISM by a type II SN exploding in a medium of average density 1 cm<sup>-3</sup> and solar metallicity was estimated by Thornton et al.(1998) to be $`7\times 10^{50}`$ ergs at the time of the peak luminosity (i.e., $`t_0`$), dropping to $`0.2\times 10^{50}`$ ergs at $`t13t_0`$, roughly equally distributed in thermal and kinetic energy. In our models, $`t_0`$ coincides with the time at which the energy is released for each SN explosion and distributed in a volume of radius $`11.7`$ pc. However, the estimates of Thornton et al. (1998) should be regarded with some caution essentially because they are based on a one-dimensional model in which the gas has few degrees of freedom. Furthermore, the cooling curve adopted by Thornton et al. (1998) excludes the radiative cooling from neutral atoms for temperatures lower than 10<sup>4</sup> K and does not include other potential cooling mechanisms such as neutrino cooling, bremsstrahlung radiation and cooling by thermal conduction. For the investigation of the supernova rate effect on the velocity dispersion, we adopt a conservative value of the feedback efficiency of $`ฯต=0.25`$.
## 5 THE SUPERNOVA RATE-VELOCITY DISPERSION RELATION
The most relevant parameter in our simulations is the supernova rate $`\eta `$. Fixing $`ฯต`$ at a value of 0.25, we performed a set of simulations with different values of $`\eta `$ ranging from 0.01 to 10 times the Galactic value. The dependence of $`v_c`$ and $`\sigma `$ on $`\eta `$ is displayed in Fig. 12 and Fig. 13, respectively. For $`\eta /\eta _G0.5`$, Fig. 12 and Fig. 13 show that $`v_c`$ and $`\sigma `$ increase rapidly with an increasing supernova rate, thus mimicking a starburst regime similar to the one observed in Fig. 2. A fit to the data for $`\eta /\eta _G`$ between 1 and 10 is $`\sigma =0.78(\pm 0.26)\eta /\eta _G+8.30(\pm 1.36)`$ km s<sup>-1</sup>. This relation might be useful for semi-analytical modeling of the central regions in galaxies, e.g., the modeling of active galactic nuclei (AGNs), in which hydrostatic equilibrium requires that the effective gravitational potential be balanced by the turbulent pressure of the gas. In current analytical models of AGNs, the intrinsic turbulent velocity of the gas is neglected and gas clouds are assumed to have the same velocity dispersion of the nuclear stellar cluster (e.g., Schartmann et al. 2005). For $`\eta /\eta _G0.5`$, the velocity dispersion shows a slow decrease with a decreasing $`\eta `$. This transition at $`\eta /\eta _G0.5`$ could indicate that the velocity dispersion of the medium in the low rate regime is not fixed by SN driving alone. SNe explosions will cause a certain fraction of the gas to be maintained in the thermally unstable regime when cold gas is restored to the warm phase. Thermal instability occurs in the cooling expanding shells, but also everywhere in the inter-supernova remnants medium where the criterion for thermal instability (see Eq. 9 in Vรกzquez-Semadeni, Gazol & Scalo 2000) is satisfied, thus adding an extra component to the kinetic energy injected into the medium. Since we are using the same cooling curve in the low temperature regime as Vรกzquez-Semadeni et al. (2000), the fit used by the latter authors and by Spitzer (1978) for the data of Dalgarno & McRay (1972) remains valid and the thermally unstable regime will be confined in the temperature range 398 K $``$ $`T`$ $``$ 10000 K.
We interpret the flatness of the $`\sigma \eta `$ relation for $`\eta /\eta _G0.5`$ as resulting from the interplay between direct supernova driving and thermal instability. Fig. 14 shows, after convergence is reached, the dependence of the volume filling factor of the thermally-unstable gas on the normalized SN rate, whereas Fig. 15 and Fig. 16 display, as examples, the time evolution of the volume filling factor of the unstable (398 K $`<T<`$ 10000 K), cold ($`T<398`$ K) and warm gas ($`T>10000`$ K) in the models with ($`\eta /\eta _G,ฯต`$)=($`0.01,0.25`$) and ($`\eta /\eta _G,ฯต`$)=($`0.1,0.25`$), respectively. In Fig. 14, the volume filling factor of the unstable gas which is an indicator of the occurence of TI shows a transition at $`\eta /\eta _G0.5`$ which roughly corresponds to the position of the transition observed in the $`\eta v_c`$ and $`\eta \sigma `$ relations, and a non-zero value at smaller values of $`\eta `$. For the large SN rate values, the gas is predominantly hot, becoming increasingly hotter with time. Therefore the volume filling factor of the unstable gas in that regime is close to zero. In the regime where $`\eta /\eta _G0.5`$, TI is more efficient in converting the gas into the cold phase at the lower end values (i.e., TI is more efficient at $`\eta /\eta _G=0.01`$ than at $`\eta /\eta _G=0.1`$). Fig. 15 and Fig. 16 show that the converged value of the volume filling factor of the cold gas in the simulation with ($`\eta /\eta _G,ฯต`$)=($`0.01,0.25`$) is $`F_c0.8`$, whereas this value is $`F_c0.5`$ for the simulation with ($`\eta /\eta _G,ฯต`$)=($`0.1,0.25`$). For the same average density in both simulations, SNe exploding in a medium with larger fractions of its volume in the cold phase as in the simulation with ($`\eta /\eta _G,ฯต`$)=($`0.01,0.25`$) will evolve in a cooled lower pressure environment than in the simulation with ($`\eta /\eta _G,ฯต`$)=($`0.1,0.25`$). Numerical simulations of SN driven turbulence by Kim (2004) show that the evolution of SNe remnants in media with lower external pressures leads to higher velocity dispersions of the gas. The existence of a background heating process could modify to some extent the present conclusion. However, the background heating should be strong enough to oppose the dramatic cooling of the medium in the low SN regime and help maintain large fractions of the gas at higher temperatures. However, in the absence of an external heating mechanism to the galaxy, the backgournd heating can only be a small fraction of the energy released by SN explosions and might not play a significant role (see also the discussion in ยง 8). Finally, it is clear that for very small SN rates, SN driving would not be able to sustain any turbulence in the medium. Turbulence will then decay before the next SN explosion occurs.
## 6 THE EFFECT OF THE AVERAGE DENSITY
As the simulations described in this chapter are not scale free, because of the presence of a realistic cooling function, another relevant point to investigate is the role of the average density. Empirical star formation laws state that the star formation rate decreases with decreasing gas surface density (Schmidt 1959,1963; Kennicutt 1998a,1998b; Dopita & Ryder 1994; Prantzos & Silk 1998). In all previous simulations, we have used an average density of 0.5 cm<sup>-3</sup>. It is particularly interesting to test the effect of varying the density for small values of the supernova rate. Over-plotted on Fig. 13 are the results of two simulations with $`\eta /\eta _G=0.05`$ and $`\eta /\eta _G=0.01`$ where the average density have been decreased, by a factor of 5 and 10, respectively. These are shown with the full square and full hexagon, respectively. Dropping the scaling coefficient, if one assumes the supernova rate (i.e., star formation rate)-gas density to follow a Kennicutt type law $`\eta /\eta _G=\overline{n}^{1.4\pm 0.15}`$ (Kennicutt 1998a,1998b), for an average density value $`\overline{n}=0.5`$ cm<sup>-3</sup> corresponds $`\eta /\eta _G=0.38\pm _{0.037}^{0.042}`$. To values of $`\eta /\eta _G=0.05`$ and 0.01 will correspond, using the same law, average densities of $`\overline{n}=0.117\pm _{0.026}^{0.027}`$ cm<sup>-3</sup> and $`\overline{n}=0.037\pm _{0.0123}^{0.0139}`$ cm<sup>-3</sup>, which roughly equal the densities of 0.1 cm<sup>-1</sup> and 0.05 cm<sup>-1</sup> we have adopted for those rates. Thus, we can consider the three points ($`\eta /\eta _G,\overline{n}`$\[cm$`{}_{}{}^{3}]`$)=$`(0.5,0.5)`$, (0.05,0.1) and (0.01,0.05) as being a rough representation of a Kennicutt law in the $`\eta \sigma `$ space. Tentatively, the preliminary conclusion we can draw here is that the flatness of the $`\sigma \eta `$ relation around $`6`$ km s<sup>-1</sup> can be maintained at low values of the supernova rate if the average density is reduced for lower rates, as predicted by the Kennicutt law. However, more simulations are needed to confirm this result and to probe the results for other star formation laws. Unfortunately, this is beyond the scope of this work, essentially for reasons of CPU time.
## 7 COMPARISON TO THE OBSERVATIONS
In Fig. 17, the same simulations appearing in Fig. 13 are shown after transforming the SN rate per unit volume into a star formation rate per unit area (units of M yr<sup>-1</sup> kpc<sup>-2</sup>). We use the transformation of the SN rate into a star formation rate (SFR), $`\eta `$/SFR=0.0067, derived using the PEGASE stellar population synthesis model (Fioc & Rocca-Volmerange 1997) and assuming a Salpeter Initial Mass Function (IMF) (Salpeter 1955). We perform a comparison to two galaxies, NGC 628 and NGC 6946 for which the star formation rates have been estimated, at different radii, from H$`\alpha `$ observations (Martin & Kennicutt 2001), along with velocity dispersion estimates which are derived from H i 21 cm line observations (Shostak & van der Kruit 1984 and Kamphuis & Sancisi 1993). Fig. 18 shows a comparison of the velocity dispersion measured from the H i gas (100 K $`T`$ 12000 K) velocity profile $`\sigma _{\mathrm{H}\mathrm{i}}`$ to the same observations. A number of remarks can be drawn from the comparisons presented in Fig. 17 and Fig. 18 : (a) The position of the transition to the starburst regime (i.e., location of sharp increase in the velocity dispersion) at around SFR/Area $`5\times 10^310^2`$ M yr<sup>-1</sup> kpc<sup>-2</sup> observed in the simulations, is relatively in good agreement with the observations in NGC 628 and NGC 6949. It is also in very good agreement with the transition to the starburst regime observed in Fig. 2, (b) there is good agreement between our models and the observations at the high SN rate values within $`23`$ km s<sup>-1</sup>. This difference can be easily explained by the effect of beam smearing which tends to increase the observed velocity dispersions, particularly in the inner parts of galaxies (c) At intermediate and low SN rates, $`\sigma `$ and $`\sigma _{\mathrm{H}\mathrm{i}}`$ fall below the observed values by a factor of 2-3 even when the density correction related to the Kennicutt law is taken into account in the case of $`\sigma _{\mathrm{H}\mathrm{i}}`$ (open square and star in Fig. 18). The values of $`\sigma `$ and $`\sigma _{\mathrm{H}\mathrm{i}}`$ become very similar at the low rates as most of the gas in these simulations has temperatures that are below $`12000`$ K (see the profiles in Fig. 9 and Fig. 10). On the other hand no H i gas is found in the simulations with $`\eta /\eta _G1`$. It should however be kept in mind that we have adopted a rather conservative value for the supernova feedback efficiency (i.e., $`ฯต=0.25`$). Fig. 11 shows that for $`ฯต=0.5`$ and 1, for the SN rate value of $`\eta /\eta _G=0.1`$, the velocity dispersions of the H i gas are of the order $`5`$ and $`6`$ km s<sup>-1</sup>, respectively, in closer agreement with the observations.
On the other hand, the velocity dispersions observed in NGC 628 and NGC 6946 (same for the observations presented in Fig. 2), in addition of being affected to some degree by beam smearing due to the limited resolution of radio telescopes (this is a minor effect at the outer galactic radii), are global velocity dispersions which do not disentangle the true dynamical value, which is the only one we have in our simulations, from the thermal broadening of the line, simply because the temperature structure of the gas in the observations is not known. Thermal broadening, $`v_T=(2k_bT/m)^{1/2}`$, is the dispersion of the velocity probability distribution function for an ensemble of particles of a non relativistic and non degenerate gas which is in thermal equilibrium at a kinetic temperature $`T`$, where $`m`$ is the mass of the particles and in our case, $`m`$ is the equal to the proton mass. Seeking a better match between our models and the observations, we have corrected for the effect of thermal broadening in two ways :
(a) the simplistic method : We can subtract, quadratically, from the observational velocity dispersion values, a velocity component associated with thermal broadening at a given equilibrium temperature for the H i gas, $`\sigma _{dyn}=(\sigma _{tot}^2v_T^2)^{1/2}`$, where $`\sigma _{dyn}`$ is the dynamical component that is compared with the dynamical component observed in the simulations and $`\sigma _{obs}`$ and $`v_T`$ are the observed velocity dispersion and the above mentioned thermal component, respectively. This is obviously a simplification as we assume the H i gas to have the same temperature at all radii. Fig. 19 shows the corrected velocity dispersions, using this method, for three values of the equilibrium temperature of the H i gas, namely 100 K, 500 K, and 2000 K. The correction improves the agreement between the observations and the simulations, particularly at the outer galactic radii and if the H i is assumed to be warmer than 100 K (bottom plots in Fig. 19 where the H i gas in NGC 628 and NGC 6946 is assumed to be at 2000 K). The fact that the gas in the outer parts could have temperatures higher than 100 K is plausible since the density in the outer parts is low and the gas could be more easily heated by cosmic rays and photoelectric effect. This is unlikely however to be the case for the inner parts of the galaxy where the H i gas is more likely to have, like in the Milky Way, a non-negligible fraction of the H i in the cold phase at around 100 K. Nevertheless, recent observations by Heiles (2001) suggest that about half of the mass of the diffuse interstellar gas in the Galaxy may have temperatures which are larger than 100 K (a few hundreds to a few thousands Kelvin). Only an accurate determination of the temperature structure in galaxies such as NGC 628 and NGC 6946 may lay out strong constraints on the contribution of thermal broadening as a function of radius to the total velocity dispersion.
(b) the less simplistic method : Here, we have assumed that the particles in each cell have a Gaussian velocity profile which is centered around the local dynamical velocity and which have a dispersion in the velocity space equal to $`v_T`$, where $`v_T`$ is the local thermal broadening calculated using the local temperature. The amplitude of the profile is given by the local density. Only cells with $`T12000`$ K and $`\overline{n}0.25`$ cm<sup>-3</sup> are taken into account. The individual velocity profiles are summed up in the velocity space and binned with a spectral bin size of 1 km s<sup>-1</sup>. Fig. 20 and Fig. 21 display two mass weighted, thermally broadened velocity profiles corresponding to simulations ($`\eta /\eta _G,ฯต`$)=(1,0.25) and (0.1,0.25), respectively. They can be compared to Fig. 8 and Fig. 10, respectively, in order to appreciate the effects of thermal broadening on the line profile. Fig. 22 shows a comparison of the H i velocity dispersion $`\sigma _{\mathrm{H}\mathrm{i}}`$ to the observational data using this more reliable approach for correcting for the effect of thermal broadening. The result is somewhat encouraging. For SFR/Area in the range of $`5\times 10^310^2`$ M yr<sup>-1</sup> kpc<sup>-2</sup> (i.e., $`\eta /\eta _G0.51`$), the agreement to the observations concerning both NGC 628 and NGC 6946 is quite acceptable. At the lower SN rates values, the agreement is less satisfying even when the Kennicutt rate-adapted average densities are used (open square and star in Fig. 22). However, a comparison only to the data of NGC 628 and NGC 6946 might be slightly misleading. In the case of galaxies such as NGC 1058 (Dickey, Hanson & Helou 1990) and NGC 3938 (van der Kruit & Shostak 1982), for which we unfortunately do not have radially dependent estimates of the star formation rate, the velocity dispersion levels off in the outer radii at a value of the order of 5-6 km s<sup>-1</sup> which is in better agreement to the values coming out from our simulations.
## 8 THE NEED FOR IMPROVED MODELS
Effect of the background heating : An important effect, which is not accounted for in our models is a background heating of the gas by the photoelectric effect, cosmic rays and soft X-rays (see also ยง2). These background heating processes might be of little importance in the case of media with SFR/Area $`5\times 10^3`$ M yr<sup>-1</sup> kpc<sup>-2</sup> (i.e., $`\eta /\eta _G0.5`$), but might play a significant role in maintaining a warmer phase of the gas at the lower SN rates. Thus, the velocity line profiles might be broader than what we have calculated in the absence of such processes. Kim (2004) showed that the velocity dispersion of the gas in a SN driven medium is reduced if the average pressure of the gas in increased. In the absence of background heating, the gas cools efficiently to the minimum temperature of 100 K in the regions where SN explosions are rare (i.e., this is particularly true for the simulations with the low SN rates). Hence, the gas pressure is reduced and the velocity dispersion enhanced. On the other hand, the shock produced by a SN explosion expanding in a cold medium would lead to higher compressions than a shock propagating in a warmer medium. The stronger compressions would lead to an enhanced cooling in the compressed SN shells which causes the energetic content of the explosion to be depleted faster. The intensity of the background heating is difficult to estimate for systems with different SN rates. As no established formulations of this problem exist in the literature, we intend, in future work, to model the background heating as being a fraction of the total SN heating and quantify its effects on the resulting velocity dispersion of the gas.
Effect of the vertical structure : The vertical stratification, which we have neglected in this work, might be one of the physical effects that we need to include first in subsequent models. We intend to perform models with a much larger length of the box in the vertical direction, using outflow boundary conditions to allow the gas to escape from the upper and lower boundaries of the simulation box. The escaping hot gas would not affect the H i 21 cm line profile. However, denser blobs of gas which are expelled at higher latitudes by the SN explosions will cool and fall back into the disk in the form of high velocity clouds. In principle, these high velocity clouds should be observed as broad wings in the velocity profile and are usually fitted with a second, broader Gaussian function. However, in-falling decelerated gas, close to the galactic disk will be most probably mixed with the local gas in the disk, thereby contributing, eventually, to the broadening of the velocity profile. The impact of HVCs on the galactic can also substantially enhance the local level of turbulence particularly for massive HVCs.
Effect of the chemistry : In the simulations presented in this work, the cooling curve we have used is a solar metallicity curve that assumes chemical equilibrium. This cooling curve describes principally the radiative cooling by neutral atoms whereas molecular cooling by molecules such as the H<sub>2</sub> molecule is neglected. The additional cooling at lower temperatures might enhance the local effects of TI and increase the value of the velocity dispersion. At the high and intermediate SN rates, this might not be of much relevance as the over-densities that are produced are a factor 3-5 the average density (a few cm<sup>-3</sup>) (see Fig. 5) and no dense molecular material is expected to form. At low rates, clouds have time to form and condense further before being destroyed by the next generation of SNe (see Fig. 6). Molecular hydrogen, starts to form when densities of the order of $`10^3`$ cm<sup>-3</sup> are reached and which become shielded against UV radiation (Bergin et al. 2004). Such densities are not reached in our simulations, essentially because of the limitations in the numerical resolution. The introduction of a simple chemical network to follow locally the fraction of molecular hydrogen in a hydrodynamical adaptative-mesh-refinement (AMR) code would help tackle the problem of molecules formation in the expanding shells more accurately, and help better account for the additional molecular cooling and its effects on the velocity field.
Effect of the metallicity : The occurrence and efficiency of the thermal instability is intimately related to the shape of the cooling curve which characterizes the medium. Cooling curves are a reflection of the strength of the emission lines of atoms present in the medium. At lower metallicity, the emission lines from metals are weaker and the cooling less efficient. The dependence of the cooling rate on the metallicity has been calculated by Boehringer & Hensler (1989), unfortunately only in the temperature range $`10^410^8`$ K. Their results show that the cooling rates may differ, in some temperature ranges, by several orders of magnitude for metallicities between $`10^2Z/Z_{}2`$ ($`Z/Z_{}=2`$ is the upper metallicity limit in their calculations) and becomes independent of the metallicity for $`Z/Z_{}<10^2`$. In future work, we plan to investigate the role of metallicity on the dynamics of the ISM by including it as a parameter in our simulations.
Effect of self-gravity : Kim et al. (2003) in some of their MRI simulations which include gravity show that for a single-phase medium and for a Toomre parameter $`Q_{th}1.7`$ which is more appropriate for the external regions of galactic disks where the gas surface density is low, self-gravity can be responsible for only 20 $`\%`$ of the level of $`1.63.2`$ km s<sup>-1</sup> turbulence generated by the MRI. However, in our models with SN explosions, the presence of self-gravity, provided enough numerical resolution is affordable to resolve the expanding shells, might lead to the development of Rayleigh-Taylor and Kelvin-Helmotz instabilities (Shu 1992), particularly in the case of SN remnants evolving in non-spherical environments (see e.g., Gazol-Patiรฑo & Passot 1999).
Effect of the magnetic field : In the present work, magnetic fields have been neglected. Several authors have performed simulations of a supernova explosions in a magnetized medium (e.g., Ferriรจrre 1991; Gazol-Patiรฑo & Passot 1999; Korpi et al. 1999a,b; Kim 2004; de Avillez & Breitschwerdt 2005; Mac Low et al. 2005). The latter simulations agree that the effect of magnetic fields in essentially to oppose the radial expansion of a supernova remnant, thus reducing the energy transmitted to the ISM and reducing the velocity dispersion. In particular, the results of Kim (2004) show that the velocity dispersion in a SN driven medium is related to the total (thermal + magnetic) pressure. Dispersions are smaller in environments with higher total pressures (i.e., higher magnetic field values for a given thermal pressure). For a comparison, at the Galactic SN rate we find that the one-dimensional total gas velocity dispersion is $`8.5`$ km s<sup>-1</sup>, whereas in the magnetized case Kim (2004) finds $`9`$ and $`10`$ km s<sup>-1</sup> for a weak magnetic field value of 2 $`\mu `$G associated to an average density of 0.2 cm<sup>-3</sup> in the directions parallel to the mean field and perpendicular to it, respectively. For a stronger field value of 8 $`\mu `$G associated to an average density of $`0.8`$ cm<sup>-3</sup>, Kim (2004) finds velocity dispersions of $`5`$ and $`7`$ km s<sup>-1</sup> in the directions parallel and perpendicular to the field, respectively. In view of Kimโs (2004) results, we speculate that in our simulations, as SNe occur in region of higher density, the existence of a magnetic field which would be compressed in those regions, would have the effect of lowering our measured values of the total gas velocity dispersion by a factor of a few tens of percent (typically 10-20 $`\%`$).
## 9 SUMMARY AND DISCUSSION
In this paper, we investigated the dependence of the velocity dispersion in the interstellar medium (ISM) on the supernova (SN) rate $`\eta `$, the SN feedback efficiency $`ฯต`$, and in some cases on the ISM average number density $`\overline{n}`$. We use local, three-dimensional numerical simulations in which SN type II explosions are detonated in random positions of the grid, separated by time intervals which are inversely proportional to the SN rate. Radiative cooling of the gas is also taken into account with a minimum cutoff temperature of 100 K. For the purpose of simplifying the problem, other physical processes and characteristics of galactic disks such as the vertical stratification, magnetic fields and gravity are neglected. For each simulation, we calculate the three-dimensional characteristic velocity dispersion $`v_c`$ (Eq. 1) and the one-dimensional velocity dispersion, $`\sigma `$, obtained by fitting the line of sight velocity profile with a Gaussian function. We also calculate $`\sigma _{\mathrm{H}\mathrm{i}}`$, which is the one-dimensional velocity dispersion obtained from a Gaussian fit of the line of sight velocity profile of the gas with a temperature $`T12000`$ K and a number density $`n0.25`$ cm<sup>-3</sup> (i.e., the H i gas). Ideally, a full investigation of the two-dimensional parameter space ($`\eta ,ฯต`$) would be necessary, however, this would lead us beyond our current computational capabilities. At this stage, we have taken an intermediate approach and explored the effects of $`\eta `$ and $`ฯต`$ independently by fixing one parameter and varying the other.
Our results show that $`v_c`$, $`\sigma `$ and $`\sigma _{\mathrm{H}\mathrm{i}}`$ depend on the SN feedback efficiency $`ฯต`$ as $`A\sqrt{ฯต}`$, where $`A`$ is a scaling coefficient, different for each quantity. This is expected as the velocity has the dimensions of the square root of an energy. In a second set of simulations, we fixed $`ฯต`$ to a value of 0.25 and varied the SN rate in the range \[0.01,10\] $`\eta _G`$, where $`\eta _G`$ is the Galactic SN rate. We compare the velocity dispersion for the simulations with the different SN rates to NGC 628 and NGC 6964 for which both the velocity dispersion and the the star formation rates are radially resolved. The dependence of the velocity dispersion on the SN rate is complex. For values of $`\eta /\eta _G0.5`$ (i.e., SFR/Area $`10^2`$ M yr<sup>-1</sup> kpc<sup>-2</sup>), $`v_c`$ and $`\sigma `$ increase sharply with increasing values of $`\eta `$. For $`\eta /\eta _G0.5`$, $`v_c`$ and $`\sigma `$ show a slower decrease with a decreasing $`\eta `$. This transition at $`\eta /\eta _G0.5`$ is probably an indication that the velocity dispersion of the medium in the low rate regime is not fixed by SN driving alone. We interpret the flatness of the $`\sigma \eta `$ relation as resulting from the efficient development of thermal instability (TI) in the medium with the low SN rates. We quantify the efficiency of TI by evaluating the volume filling factor of the unstable gas $`F_u`$ for the different SN rates. $`F_u`$ appears to be correlated with the velocity dispersion. For large values of the SN rate $`\eta /\eta _G0.51`$, $`F_u`$ is close to zero, with a transition to non zero values at $`\eta /\eta _G0.5`$, peaks at $`\eta /\eta _G0.1`$, before decreasing slowly at smaller SN rates. Interestingly, the position of the transition to the starburst regime (high SFR rates associated with high velocity dispersions) is in relatively good agreement with the one seen in the observations (see Fig. 2, Fig. 18 and Fig. 19).
We compare $`\sigma _{\mathrm{H}\mathrm{i}}`$ to the observations as a function of the SN rate (i.e., SFR rates). $`\sigma `$ and $`\sigma _{\mathrm{H}\mathrm{i}}`$ have nearly similar values at the low SN rates as most of the gas in this regime is H i gas. H i gas is not present in the simulations with $`\eta /\eta _G1`$. $`\sigma _{\mathrm{H}\mathrm{i}}`$ is observed to be nearly independent of the SN rate, leveling off at $`3`$ km<sup>-1</sup>. This value is a factor $`23`$ lower that the velocity dispersion plateau observed in galaxies such as NGC 1058, NGC 628 and NGC 6964 (Fig. 1 and Fig. 18). However, the set of simulations with $`\eta /\eta _G=0.1`$ and the different feedback efficiencies (i.e, Fig. 11) suggests that H i gas velocity dispersions of the order of 5-6 km s<sup>-1</sup> can be obtained for feedback efficiencies $`ฯต0.5`$, in closer agreement with the observations. For the simulations with our fiducial value of $`ฯต`$, we have corrected the H i line velocity profiles by accounting for the effect of thermal broadening. A reasonable agreement is found between the models and the observations for values of the SFR/Area in the range $`5\times 10^310^2`$ M yr<sup>-1</sup> kpc<sup>-2</sup>. For smaller values of the SFR, the fact that the velocity dispersions are a factor $`2`$ smaller than the observed values could result from the fact that we have underestimated the SN feedback efficiency. Otherwise, it might be an indication of the existence of secondary heating and/or driving mechanisms in the outer parts of galaxies where the star formation rate is low. An investigation of the effects of SN driving in the presence of other physical processes and gas instabilities is essential and is left to future work.
It is a pleasure to thank Enrique Vรกzquez-Semadeni, Fabian Walter, Hans-Walter Rix, Axel Brandenburg, Maarit Korpi, Joop Schaye and Robert Piontek for useful comments and discussion. We are also very grateful to the anonymous referee for interesting comments and suggestions. S. D. would like to thank Hans-Walter Rix and Thomas Henning for the financial support at the MPIA. Calculations have been performed on the MPIAโs SGI Origin 2000 located at the Rechnenzentrum of the Max-Planck Gesellschaft, Garching. ZEUS-3D was used by courtesy of the Laboratory of Computational Astrophysics at the NCSA. This research has made use of NASAโs Astrophysics Data System Bibliographic Services.
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# Distribution of current in non-equilibrium diffusive systems and phase transitions
## I Introduction
Stochastic lattice gas models have been extensively studied recently as they are among the simplest examples of non-equilibrium systems. A powerful approach to understand their steady state was developed by Bertini, De Sole, Gabrielli, Jona-Lasinio, Landim, as a macroscopic fluctuation theory (MFT) which gives, for large diffusive systems, the probability distribution of trajectories in the space of density profiles bdgjl ; bdgjl2 . The MFT relies on the hydrodynamic large deviation theory kov ; spohn ; KL which provides estimates for the probability of observing atypical space/time density profiles. It gives a framework to calculate a large number of properties of stochastic lattice gas, such as the large deviation functional of the density profiles. Recent developments of the hydrodynamic large deviation theory bdgjl3 ; bd enabled to estimate also the large deviations of the current through the system. What the MFT provide are the equations of the time evolution of the most likely density profile responsible of a given fluctuation. What it does not provide, in general, is the solution of these equations which would give quantitative predictions for the distribution of the fluctuations. So far, for the large deviation functional of the density profiles in the steady state, the equations could only be solved in a few cases of non-equilibrium systems with open boundaries (the symmetric exclusion process bdgjl2 , the Kipnis, Marchioro, Presutti model bgl ; kmp ). For the SSEP the results of the MFT were in full agreement with the results obtained dls ; dls2 ; ed from the exact knowledge of the weights of the microscopic configurations in the steady state.
In our previous work bd , we developed a theory to calculate the large deviation function of the current through a long one dimensional diffusive lattice gas in contact at its two ends with two reservoirs at unequal densities. Our approach was based on an assumption, the additivity principle, which relates the large deviation function (LDF) of the current of a system to the LDFโs of subsystems, when one breaks a large system into large subsystems. This assumption is in fact equivalent to the hypothesis, within the hydrodynamic large deviation framework, that to observe, for a very long time period $`T`$, an average current $`q=Q_T/T`$, the system adopts a profile with a shape, fixed in time, but of course depending on $`q`$ (here $`Q_T`$ is the total number of particles transfered, say from the left reservoir to the system during time $`T`$). The additivity principle allows one to obtain explicit expressions bd for all the cumulants of the integrated current $`Q_T`$. The predictions of our theory were tested in a few cases bd and the results were found in complete agreement with what was already known or what could be derived by alternative approaches ddr ; wr ; hrs .
Recently, it was pointed out bdgjl3 that even if our predictions bd are valid for some diffusive lattice gas, it might happen that, to produce an average current $`q`$ over a long period of time, the best profile is time-dependent. One of the goals of the present work is to show that, for a simple example, the weakly asymmetric exclusion process on a ring, this indeed happens for some range of parameters.
Let us consider, as we shall do it in the rest of this paper, the time evolution of a one dimensional stochastic lattice gas on a lattice of $`N`$ sites. According to the hydrodynamic formalism spohn , a given lattice gas can be characterized by two functions $`D(\rho )`$ and $`\sigma (\rho )`$ of its density $`\rho `$. One way to define them bd is to consider a one dimensional system of length $`N`$ connected to reservoirs at its two ends. For such a lattice gas, the variance of the total charge $`Q_T`$ transfered during a long time $`T`$ from one reservoir to the other is given, for large $`N`$, by definition of $`\sigma (\rho )`$ by
$$\frac{Q_T^2}{T}=\frac{\sigma (\rho )}{N}$$
(1)
when both reservoirs are at the same density $`\rho `$. On the other hand if the left reservoir is at density $`\rho +\mathrm{\Delta }\rho `$ and the right reservoir at density $`\rho `$, the average current is given, for small $`\mathrm{\Delta }\rho `$, by
$$\frac{Q_T}{T}=\frac{D(\rho )\mathrm{\Delta }\rho }{N}$$
(2)
which is simply Fickโs law and defines the function $`D(\rho )`$. In the symmetric simple exclusion process $`\sigma (\rho )=\rho (1\rho )`$ and $`D(\rho )=1/2`$ spohn whereas in the Kipnis, Marchioro, Presutti model bgl ; kmp $`\sigma (\rho )=\rho ^2`$ and $`D(\rho )=1/2`$. The effect of a uniform weak electric field of strength $`\nu /(2N)`$ acting from left to right on the particles is to modify (2) into
$$\frac{Q_T}{T}=\frac{D(\rho )\mathrm{\Delta }\rho }{N}+\frac{\nu \sigma (\rho )}{N}$$
(3)
This equation follows from the linear response theory spohn .
Once $`D(\rho )`$ and $`\sigma (\rho )`$ are known for a given diffusive system, the probability of observing the evolution of a density profile $`\rho (x,s)`$ and a rescaled current $`j(x,s)`$ for $`0<s<T`$ during a time $`TN^2`$ is given, according to the hydrodynamic large deviation theory bdgjl3 , by
$$\mathrm{Pro}(j(x,s),\rho (x,s))\mathrm{exp}\left[\frac{_{[0,T]}^\nu (j,\rho )}{N}\right]$$
(4)
where $`_{[0,T]}^\nu `$ is defined by
$`_{[0,T]}^\nu (j,\rho )={\displaystyle _0^T}๐s{\displaystyle _0^1}๐x`$
$`{\displaystyle \frac{[j(x,s)+D(\rho (x,s))\rho ^{}(x,s)\nu \sigma (\rho (x,s))]^2}{2\sigma (\rho (x,s))}}`$
with $`\rho ^{}=\rho /x`$ and where the rescaled current $`j(x,s)`$ is related to the density profile $`\rho (x,s)`$ by the conservation law
$$\frac{d\rho (x,s)}{ds}=\frac{dj(x,s)}{dx}$$
(6)
A formalism equivalent to this hydrodynamic large deviation theory was developed independently pjsb ; jsp in the context of the full counting statistics of the transport of free fermions through disordered wires. A simple derivation of (I) and (6) is given in the appendix.
The large deviation function $`G(j_0)`$ of the current is then defined as
$`\mathrm{Pro}\left({\displaystyle \frac{Q_T}{T}}={\displaystyle \frac{j_0}{N}}\right)\mathrm{exp}\left({\displaystyle \frac{T}{N}}G(j_0)\right)`$ (7)
for large $`T`$ and $`N`$
(In (7), one has first to take the limit $`T\mathrm{}`$ and then make $`N`$ large; in practice (7) should hold when $`TN^2`$ as $`N^2`$ is the characteristic time of a diffusive system of size $`N`$).
Now according to (4), the large deviation function $`G(j_0)`$ is given by
$$G(j_0)=\underset{T\mathrm{}}{lim}\left[\frac{1}{T}\underset{\rho (x,s)}{\mathrm{min}}_{[0,T]}^\nu (j,\rho )\right]$$
(8)
where the current $`j(x,s)`$ satisfies for large $`T`$ and all $`x`$ the constraint
$$\underset{T\mathrm{}}{lim}\frac{1}{T}_0^Tj(x,s)๐s=j_0$$
(9)
with the profile $`\rho (x,s)`$ and the current $`j(x,s)`$ connected by (6).
In the following, we will often consider, instead of (7), the generating function of the current:
$$e^{\lambda Q_T}e^{T\mu (\lambda )}\mathrm{for}\mathrm{large}T$$
(10)
and then, according to (7,8), $`\mu (\lambda )`$ is given by
$`\mu (\lambda )={\displaystyle \frac{1}{N}}\underset{j_0}{\mathrm{max}}[\lambda j_0+G(j_0)]=`$
$`{\displaystyle \frac{1}{N}}\underset{T\mathrm{}}{lim}{\displaystyle \frac{1}{T}}\underset{\rho (x,s)}{\mathrm{max}}\left[\lambda {\displaystyle _0^T}j(x,s)๐s_{[0,T]}^\nu (j,\rho )\right]`$
For a system of length $`N`$ connected to two reservoirs at densities $`\rho _a`$ and $`\rho _b`$ at its two ends, the calculation of the large deviation function $`G(j_0)`$ of the current is therefore reduced to finding the time-dependent profile $`\rho (x,s)`$ which optimizes (8) under the constraints (6) and (9) and with the additional boundary conditions $`\rho (0,s)=\rho _a`$ and $`\rho (1,s)=\rho _b`$. All the results of our previous work bd follow then from the assumption that this optimal profile does not vary with time (except from boundary effects near time 0 and time $`T`$ which do not contribute in the large $`T`$ limit).
For a system on a ring of $`N`$ sites, as we shall consider here, the optimization problem is the same except for the boundary condition which becomes $`\rho (0,s)=\rho (1,s)`$ and the fact that the total density $`\rho _0`$ on the ring becomes an additional conserved quantity
$$_0^1\rho (x,s)๐x=\rho _0.$$
Our paper is organized as follows: in section II, we consider a general lattice gas on a ring and show under what conditions the flat profile becomes unstable. In sections II and III, we write the large deviation function of the current, when the optimal profile has a fixed shape moving at a constant velocity. In section IV, we present exact numerical results on the weakly asymmetric exclusion process for small system sizes which give evidence that for some range of parameters, consistent with the results of section III, the optimal profile is no longer flat but becomes space-time dependent (for $`\rho _0=1/2`$ it is only space dependent). In section V, we analyze the limit of a strong asymmetry and obtain a simple expression for the large deviation function of the current, under the assumption that the optimal profile becomes in this limit a step function. In section VI, we show that, even for a strongly asymmetric case such as the totally asymmetric exclusion process, one can exhibit time-dependent profiles determined by the Jensen Varadhan functional which give, for large system size, the exact large deviation function of the current previously calculated by the Bethe ansatz.
## II Conditions for the stability of a flat profile
Consider a lattice gas on a ring of $`N`$ sites with total density $`\rho _0`$. Under the assumption that the optimal profile $`\rho (x,s)`$ is flat, that is
$$\rho (x,s)=\rho _0$$
the large deviation function $`G(j_0)`$ is given (I,8) by
$$G_{\mathrm{flat}}(j_0)=\frac{[j_0\nu \sigma (\rho _0)]^2}{2\sigma (\rho _0)}$$
(12)
which, by (I), gives for $`\mu (\lambda )`$
$$\mu _{\mathrm{flat}}(\lambda )\frac{\lambda (\lambda +2\nu )\sigma (\rho _0)}{2N}.$$
(13)
A natural question is whether one could increase $`G(j_0)`$ (and $`\mu (\lambda )`$) by adding to this flat profile some small ($`ฯต1`$) space and time dependent perturbation of the form
$$j(x,s)=j_0+ฯต[j_1(x)\mathrm{cos}(\omega s)+j_2(x)\mathrm{sin}(\omega s)]$$
in which case due to (6)
$$\rho (x,s)=\rho _0+\frac{ฯต}{\omega }[j_1^{}(x)\mathrm{sin}(\omega s)+j_2^{}(x)\mathrm{cos}(\omega s)]$$
where $`j_1(x)`$ and $`j_2(x)`$ are periodic functions of period 1. The resulting expression for $`G(j_0)`$ to second order in $`ฯต`$ is
$`G(j_0)={\displaystyle \frac{(j_0\nu \sigma (\rho _0))^2}{2\sigma (\rho _0)}}+ฯต^2{\displaystyle _0^1}dx[{\displaystyle \frac{j_1^2+j_2^2}{4\sigma (\rho _0)}}`$
$`{\displaystyle \frac{(j_1^{\prime \prime 2}+j_2^{\prime \prime 2})D(\rho _0)^2}{4\omega ^2\sigma (\rho _0)}}+{\displaystyle \frac{j_0(j_1j_2^{}j_2j_1^{})\sigma ^{}(\rho _0)}{2\omega \sigma ^2(\rho _0)}}`$
$`+(j_1^2+j_2^2)({\displaystyle \frac{j_0^2\sigma ^{\prime \prime }(\rho _0)}{8\omega ^2\sigma (\rho _0)^2}}{\displaystyle \frac{\nu ^2\sigma ^{\prime \prime }(\rho _0)}{8\omega ^2}}{\displaystyle \frac{j_0^2\sigma ^2(\rho _0)}{4\omega ^2\sigma ^3(\rho _0)}})]`$
As this expression is quadratic in the currents $`j_1`$ and $`j_2`$, the various Fourier modes are not coupled. Choosing for $`j_1(x)`$ and $`j_2(x)`$
$`j_1(x)=a\mathrm{cos}(2\pi x)+b\mathrm{sin}(2\pi x)`$
$`j_2(x)=c\mathrm{cos}(2\pi x)+d\mathrm{sin}(2\pi x)`$ (14)
one gets
$`G(j_0)={\displaystyle \frac{(j_0\nu \sigma (\rho _0))^2}{2\sigma (\rho _0)}}+ฯต^2[(adbc){\displaystyle \frac{j_0\sigma ^{}(\rho _0)\pi }{\omega \sigma ^2(\rho _0)}}`$
$`(a^2+b^2+c^2+d^2)({\displaystyle \frac{1}{8\sigma (\rho _0)}}+{\displaystyle \frac{2\pi ^4D^2(\rho _0)}{\omega ^2\sigma (\rho _0)}}`$
$`+{\displaystyle \frac{\pi ^2\sigma ^2(\rho _0)j_0^2}{2\omega ^2\sigma ^3(\rho _0)}}+{\displaystyle \frac{\nu ^2\pi ^2\sigma ^{\prime \prime }(\rho _0)}{4\omega ^2}}{\displaystyle \frac{\pi ^2\sigma ^{\prime \prime }(\rho _0)j_0^2}{4\omega ^2\sigma ^2(\rho _0)}})]`$
The flat profile is stable against the perturbation (14), if this is a negative definite quadratic form in $`a,b,c,d`$. This is achieved when for all $`\omega `$
$`{\displaystyle \frac{1}{8\sigma (\rho _0)}}+{\displaystyle \frac{2\pi ^4D^2(\rho _0)}{\omega ^2\sigma (\rho _0)}}+{\displaystyle \frac{\pi ^2\sigma ^2(\rho _0)j_0^2}{2\omega ^2\sigma ^3(\rho _0)}}`$
$`+{\displaystyle \frac{\nu ^2\pi ^2\sigma ^{\prime \prime }(\rho _0)}{4\omega ^2}}{\displaystyle \frac{\pi ^2\sigma ^{\prime \prime }(\rho _0)j_0^2}{4\omega ^2\sigma ^2(\rho _0)}}>\left|{\displaystyle \frac{j_0\sigma ^{}(\rho _0)\pi }{2\omega \sigma ^2(\rho _0)}}\right|`$
The flat profile becomes therefore unstable if
$`{\displaystyle \frac{1}{8\sigma (\rho _0)}}+{\displaystyle \frac{2\pi ^4D^2(\rho _0)}{\omega ^2\sigma (\rho _0)}}+{\displaystyle \frac{\pi ^2\sigma ^2(\rho _0)j_0^2}{2\omega ^2\sigma ^3(\rho _0)}}`$ (15)
$`+{\displaystyle \frac{\nu ^2\pi ^2\sigma ^{\prime \prime }(\rho _0)}{4\omega ^2}}{\displaystyle \frac{\pi ^2\sigma ^{\prime \prime }(\rho _0)j_0^2}{4\omega ^2\sigma ^2(\rho _0)}}<\left|{\displaystyle \frac{j_0\sigma ^{}(\rho _0)\pi }{2\omega \sigma ^2(\rho _0)}}\right|`$
i.e.
$$8\pi ^2D^2(\rho _0)\sigma (\rho _0)+(\nu ^2\sigma ^2(\rho _0)j_0^2)\sigma ^{\prime \prime }(\rho _0)<0$$
(16)
which, given (I,12,13), can be rewritten as
$$4\pi ^2D^2(\rho _0)<N\mu _{\mathrm{flat}}(\lambda )\sigma ^{\prime \prime }(\rho _0)$$
(17)
When the flat profile becomes unstable, according to (15), the current takes the form
$$j(x,t)=j_0+A\mathrm{cos}2\pi \left(xx_0\frac{j_0\sigma ^{}(\rho _0)}{\sigma (\rho _0)}t\right)$$
(18)
where the amplitude $`A`$ would be determined by expanding $`G(j_0)`$ to higher order in $`ฯต`$.
One could analyze in a similar way the stability of the flat profile against other modes by choosing $`j_1(x)=a\mathrm{cos}(2\pi nx)+b\mathrm{sin}(2\pi nx)`$ and $`j_2(x)=c\mathrm{cos}(2\pi nx)+d\mathrm{sin}(2\pi nx)`$ and the threshold (16) would become
$$8\pi ^2D^2(\rho _0)n^2\sigma (\rho _0)+(\nu ^2\sigma ^2(\rho _0)j_0^2)\sigma ^{\prime \prime }(\rho _0)<0$$
This shows that that the fundamental ($`n=1`$) is the first mode to become unstable.
## III A simple time dependent profile
The form (18) suggests that beyond the instability the optimal profile is a fixed shape moving at a constant velocity $`v`$
$$\rho =g(xvt).$$
Due to conservation law (6) the current is then
$$j(x,t)=j_0v\rho _0+vg(xvt)$$
If such a profile is the optimal profile, then the variational principle (8) reduces to
$`G(j_0)=\underset{g(x),v}{\mathrm{min}}{\displaystyle _0^1}{\displaystyle \frac{dx}{2\sigma (g(x))}}[j_0v\rho _0+vg(x)`$ (19)
$`+D(g(x))g^{}(x)\nu \sigma (g(x))]^2`$
This is of the form (the term linear in $`g^{}`$ gives a null contribution due to the periodic boundary conditions)
$$G(j_0)=\underset{g(x),v}{inf}_0^1๐x[X(g)+g^2Y(g)]$$
(20)
where
$$X(g)=\frac{[j_0v\rho _0+vg\nu \sigma (g)]^2}{2\sigma (g)}$$
and
$$Y(g)=\frac{D^2(g)}{2\sigma (g)}.$$
The optimal $`v`$ in (19) is then given by
$$v=\frac{๐x\frac{(g\rho _0)(j_0\nu \sigma (g))}{\sigma (g)}}{๐x\frac{(g\rho _0)^2}{\sigma (g)}}=j_0\frac{๐x\frac{(g\rho _0)}{\sigma (g)}}{๐x\frac{(g\rho _0)^2}{\sigma (g)}}$$
(21)
this last simplification being due to the constraint $`g(x)๐x=\rho _0`$. With this constraint and for a fixed $`v`$, a variational calculation of the optimal $`g`$ in (20) shows that $`g`$ should satisfy
$$X^{}(g)2Y(g)g^{\prime \prime }g^2Y^{}(g)=C_2$$
Multiplying both sides by $`g^{}`$ allows one to integrate once so that $`g`$ satisfies
$$X(g)g^2Y(g)=C_1+C_2g$$
(22)
where $`C_1`$ and $`C_2`$ are constants (which is an extension of the equation (15) of bd to the case of the ring).
For fixed $`j_0`$, $`\rho _0`$ and $`v`$, if one denotes by $`g_1`$ and $`g_2`$ the two extrema of the profile $`g`$ (generically, the profile $`g(x)`$ is a periodic function of period 1 with a single minimum $`g_1`$ and a single maximum $`g_2`$), one can determine the constants $`C_1`$ and $`C_2`$ by (22) in terms of $`g_1`$ and $`g_2`$ (as $`X(g_1)=C_1+C_2g_1`$ and $`X(g_2)=C_1+C_2g_2`$). The differential equation (22) determines the whole profile (up to a translation on the ring) and the constants $`g_1`$ and $`g_2`$ are then fixed by the fact that
$`{\displaystyle \frac{1}{2}}={\displaystyle _{x(g_1)}^{x(g_2)}}๐x={\displaystyle _{g_1}^{g_2}}{\displaystyle \frac{dg}{g^{}}}`$
$`={\displaystyle _{g_1}^{g_2}}\sqrt{{\displaystyle \frac{Y(g)}{X(g)C_1C_2g}}}๐g`$
and
$`{\displaystyle \frac{\rho _0}{2}}={\displaystyle _{g_1}^{g_2}}g\sqrt{{\displaystyle \frac{Y(g)}{X(g)C_1C_2g}}}๐g`$
## IV Exact numerics for the weakly asymmetric exclusion process on a ring
We wrote a program to calculate exactly $`\mu (\lambda )`$ for the weakly asymmetric exclusion process (WASEP) on a ring of $`N`$ sites with $`P=N\rho `$ particles. In the simple symmetric exclusion process (SSEP), each particle jumps to its right at rate $`\frac{1}{2}`$ and to its left at rate $`\frac{1}{2}`$ and the functions $`D(\rho )`$ and $`\sigma (\rho )`$ are given spohn by
$$D_{\mathrm{SSEP}}=\frac{1}{2}$$
(23)
$$\sigma _{\mathrm{SSEP}}=\rho (1\rho ).$$
(24)
If one introduces a weak electric field to the right, the model becomes the WASEP and the rates become $`\frac{1}{2}+\frac{\nu }{2N}`$ to the right and $`\frac{1}{2}\frac{\nu }{2N}`$.
As the evolution is a Markov process, one can build, as explained in dl ; rdd from the Markov matrix, a $`\lambda `$-dependent matrix, the largest eigenvalue of which is $`\mu (\lambda )`$ defined by (10). According to the linear stability analysis, the flat profile becomes unstable (16) for
$$j_0^2<\rho (1\rho )[\nu ^2\rho (1\rho )\pi ^2]$$
or by (17) for
$$N\mu _{\mathrm{flat}}(\lambda )<\frac{\pi ^2}{2}.$$
(25)
We have calculated the exact eigenvalue $`\mu (\lambda )`$ for lattice sizes from $`N=6`$ to $`22`$, at density $`1/2`$ for an asymmetry $`\nu =10`$ (in order to avoid negative rates for small system sizes, we have replaced in our programs the rates $`\frac{1}{2}\pm \frac{\nu }{2N}`$ by $`\mathrm{exp}[\pm \nu /N]/2`$). The results show a rather quick convergence with increasing $`N`$ towards the value corresponding to a shape determined by (19). Clearly the flat profile gives a value too low, incompatible with the numerical data.
## V A bridge between a weak and a strong asymmetry
Assuming that the optimal profile is the one discussed in section III, if one tries to make $`\nu `$ large, and one writes
$$j_0=\nu i_0$$
(26)
the moving profile which satisfies (22) becomes very steep in the regions where it varies, and it takes the form of a step function with two constant values $`g_1`$ and $`g_2`$ separated by two discontinuities
$`g(x)=g_1\mathrm{for}0<x<y`$
$`g(x)=g_2\mathrm{for}y<x<1`$ (27)
so that the parameters $`g_1,g_2`$ and $`y`$ are related to $`\rho _0`$ and $`i_0`$ by
$$\rho _0=yg_1+(1y)g_2$$
(28)
$$i_0=y\sigma (g_1)+(1y)\sigma (g_2).$$
(29)
The expression of the velocity (21) then becomes
$$v=\nu \frac{\sigma (g_2)\sigma (g_1)}{g_2g_1}.$$
(30)
Then using the fact that $`g^{}`$ vanishes when $`g(x)=g_1`$ or $`g_2`$ and replacing (26,28,29,30) into (22) implies that the constants $`C_1,C_2`$ (in (22)) vanish at order $`\nu ^2`$. Thus asymptotically in $`\nu `$, one can rewrite (22) as
$`g^2=\nu ^2{\displaystyle \frac{1}{D(g)^2(g_1g_2)^2}}[g_2(\sigma (g_1)\sigma (g))`$
$`+g_1(\sigma (g)\sigma (g_2))+g(\sigma (g_2)\sigma (g_1))]^2`$
and
$`G(j_0)=\nu ^2{\displaystyle }{\displaystyle \frac{dx}{\sigma (g)(g_1g_2)^2}}[g_2(\sigma (g_1)\sigma (g))`$ (32)
$`+g_1(\sigma (g)\sigma (g_2))+g(\sigma (g_2)\sigma (g_1))]^2.`$
If one replaces $`g`$ by its expression (27), one gets that the order $`\nu ^2`$ of $`G(j_0)`$ vanishes. The next order in the large $`\nu `$ expansion is dominated by the rounding-off of the discontinuities in (27) as given by (V). As the profile $`g(x)`$ is composed of two monotonic parts one can then use (V) into (32) and obtain for $`g_2>g_1`$
$`G(j_0)=2\nu {\displaystyle _{g_1}^{g_2}}๐g{\displaystyle \frac{D(g)}{(g_2g_1)\sigma (g)}}|\sigma (g)(g_2g_1)`$
$`\sigma (g_1)(g_2g)\sigma (g_2)(gg_1)|`$
It is remarkable that $`y`$ is not present in this expression. In the case of the weakly asymmetric exclusion process on a ring, expressions (23,24) of $`D(\rho )`$ and $`\sigma (\rho )`$ lead for large $`\nu `$ to
$`G(j_0)=\nu [g_2g_1g_1g_2\mathrm{ln}{\displaystyle \frac{g_2}{g_1}}`$
$`(1g_1)(1g_2)\mathrm{ln}\left({\displaystyle \frac{1g_1}{1g_2}}\right)]`$
In this case (30) becomes $`v=\nu (1g_1g_2)`$.
If one takes formally $`\nu =N`$, the hopping rates $`1/2\pm \nu /2N`$ become $`1`$ and $`0`$, so that the model reduces to the totally asymmetric exclusion process and one gets from (7,26,V)
$`\mathrm{Pro}({\displaystyle \frac{Q_T}{T}}=i_0)\mathrm{exp}(T[g_2g_1g_1g_2\mathrm{ln}{\displaystyle \frac{g_2}{g_1}}`$ (34)
$`(1g_1)(1g_2)\mathrm{ln}\left({\displaystyle \frac{1g_1}{1g_2}}\right)]).`$
As we will see it in section VI, this is exactly the large deviation function predicted by the Jensen-Varadhan theory JV to maintain a profile (27) formed of a shock and an antishock in the totally asymmetric exclusion process. Other aspects of the relation between the large deviation functional of the weakly asymmetric exclusion process and the Jensen-Varadhan functional in systems with open boundary conditions will be presented in bd2 .
## VI Large deviations of the current in the totally asymmetric exclusion process
In the totally asymmetric process, each particle jumps to its neighboring site, on its right, at rate 1, if the target site is empty (and there is no other jump).
The large deviation function of the current of the totally asymmetric exclusion process on a ring of $`N`$ sites, with $`P`$ particles, has been calculated exactly dl ; da . If $`Q_T`$ is the total number of jumps during time $`T`$ over a given bond on the ring, one knows that for large $`T`$,
$$e^{\lambda Q_T}e^{\mu (\lambda )T}$$
(35)
and explicit expressions of $`\mu (\lambda )`$ has been obtained for all $`N`$ and $`P`$ by the Bethe ansatz dl ; da .
For large $`N`$, it was shown in particular (equation (53) of da with the proper redefinition of the parameters) that for $`\lambda <0`$
$$\mu (\lambda )=\frac{(1e^{\lambda \rho _0})(1e^{\lambda (1\rho _0)})}{(1e^\lambda )}.$$
(36)
We are going now to argue that this result can be understood, by assuming that (35) is dominated by configurations of the form (27) moving at a velocity $`v=1g_1g_2`$. These density profiles are everywhere constant except for a shock (at some position $`z`$ with $`g(z0)=g_1`$ and $`g(z+0)=g_2`$ for $`g_2>g_1`$) and an antishock (at position $`z+y`$ with $`g(z+y0)=g_2`$ and $`g(z+y+0)=g_1`$). From the Jensen-Varadhan theory JV , the probability of maintaining such a shape moving at this velocity on the ring over a very long period of time $`T`$ reduces to the probability of maintaining an antishock between the densities $`g_2`$ and $`g_1`$ moving at velocity $`v`$. The probability of the latter event, which we denote by $`P_T(g_1,g_2)`$ is given JV by
$`P_T(g_1,g_2)\mathrm{exp}(T[g_2g_1g_1g_2\mathrm{ln}{\displaystyle \frac{g_2}{g_1}}`$ (37)
$`(1g_1)(1g_2)\mathrm{ln}\left({\displaystyle \frac{1g_1}{1g_2}}\right)]).`$
The corresponding integrated current $`Q_T`$ is
$$Q_T=T[yg_1(1g_1)+(1y)g_2(1g_2)]$$
since over a long period of time $`T`$ a given bond spends a fraction $`y`$ of the time at density $`g_1`$ and $`1y`$ at density $`g_2`$.
Therefore, if the configurations of the form (27) dominate the large deviations of the current, one expects
$`\mu (\lambda )=\underset{y,g_1,g_2}{\mathrm{max}}\{\lambda [yg_1(1g_1)+(1y)g_2(1g_2)]`$
$`[g_2g_1g_1g_2\mathrm{ln}{\displaystyle \frac{g_2}{g_1}}(1g_1)(1g_2)\mathrm{ln}\left({\displaystyle \frac{1g_1}{1g_2}}\right)]\}`$
where the maximum has to satisfy the constraint
$$\rho _0=yg_1+(1y)g_2$$
(39)
A calculation of the optimum in (LABEL:muasep), with the constraint (39) leads to
$$g_1=\frac{e^\lambda e^{\lambda (1\rho _0)}}{e^\lambda 1};g_2=\frac{e^{\lambda \rho _0}1}{e^\lambda 1}$$
and (LABEL:muasep) becomes (36).
This shows that the result of the Bethe ansatz (36) can be physically understood in terms of an optimal profile which takes the form of the step function (27). The probability of maintaining this profile is given by the Jensen-Varadhan expression (37) which in fact is identical to (34) obtained, in the large $`\nu `$ limit, for the WASEP from the hydrodynamic large deviation theory.
That the fluctuations are due, in the strong asymmetric case, to configurations formed by a gas of shocks and antishocks has been already pointed out by Fogedby f1 ; f2 ; fb . The calculation of this section shows that the large deviation of the current, in the range $`\lambda <0`$, can be understood quantitatively in terms of a single pair of shock-antishock. Whether the Fogedby theory would allow to understand all the current fluctuations, including the range $`\lambda >0`$ where the expression of $`\mu (\lambda )`$ is more complicated da than (36), remains an interesting open question.
## VII Conclusion
In the present work we have determined the limit of stability (16,17) of a flat profile for a diffusive lattice gas on a ring. This instability beyond which the optimal profile becomes modulated is of the same nature as the phase transition found for several other non-equilibrium systems bgl ; zia ; ff ; cde . As the calculation is based on a local stability analysis, one cannot of course exclude first order transitions, i.e. that the flat profile might become globally unstable.
In section IV, we have obtained numerical evidence that the macroscopic fluctuation theory predicts correctly the large deviation function of the current for the weakly asymmetric exclusion process. The numerical results are consistent with the second order phase transition predicted in section II, and with a modulated density profile moving at a constant velocity as suggested in section III. These results could in principle be confirmed by solving the Bethe ansatz equations for the WASEP, since $`\mu (\lambda )`$ can be calculated exactly for the ring geometry lk ; k .
It would be interesting to extend the results of the present work to the case of open boundary conditions. One difficulty is that the time-independent profile, found in bd , is much more complicated than the flat profile for the ring geometry, and we did not succeed so far to obtain the condition which would generalize (16,17) for this open geometry.
Lastly, we noticed that the large deviation function (34) obtained for the weakly asymmetric diffusion process in the large drift limit is identical to the one predicted for a strong asymmetry by the Jensen-Varadhan theory (37) (see also bd2 ).
Despite this bridge between the large deviation function of the current of weakly and strongly asymmetric systems, and some recent results on zero-range processes hrs , a theory of current fluctuations for strongly asymmetric lattice gas such as the ASEP with open boundary conditions remains an open problem.
## VIII Appendix: a derivation of (4,I)
We present here an heuristic derivation of the hydrodynamic large deviations (4). Let us consider a system of $`N`$ sites and decompose it into $`N/l`$ boxes of $`l`$ sites each. Let us define the density $`\rho _i(t)`$ in box $`i`$ at time $`t`$ and $`q_i(t)`$ the total number of particles transfered from box $`i`$ to box $`i+1`$ during a time interval $`t,t+\tau `$ (this time $`\tau `$ should be large enough for the $`q_i(t)`$ to be a Gaussian characterized by its average and its variance as in (1,3), but short enough compared to the characteristic time of variation of the densities $`\rho _i(t)`$).
If one writes that the $`q_i(t)`$ are Gaussian, one gets
$`\mathrm{Pro}(q_i(s),\rho _i(s))\mathrm{exp}[{\displaystyle \underset{k=1}{\overset{T/\tau }{}}}{\displaystyle \underset{i=1}{\overset{N/l}{}}}{\displaystyle \frac{1}{\frac{2\sigma (\rho _i(s))\tau }{l}}}(q_i(s)+`$
$`D(\rho _i(s)){\displaystyle \frac{\rho _{i+1}(s)\rho _i(s)}{l}}\tau \nu {\displaystyle \frac{\sigma (\rho _i(s))}{N}}\tau )^2]`$
where $`k=s\tau `$. The factor $`\nu /N`$ comes from the weak asymmetry of the jumps. Clearly the conservation of the number of particles gives
$$\rho _i(s+\tau )=\rho _i(s)+\frac{q_{i1}(s)q_i(s)}{l}$$
(40)
Now if one takes a continuous limit by writing
$$\rho _i(s)=\rho (i\frac{l}{N},\frac{s}{N^2})$$
(41)
and one defines a rescaled current by
$$q_i(s)=\frac{\tau }{N}j(i\frac{l}{N},\frac{s}{N^2})$$
(42)
one gets
$`\mathrm{Pro}(j(x,s),\rho (x,s))\mathrm{exp}[N^1{\displaystyle _0^T}ds{\displaystyle _0^1}dx`$
$`{\displaystyle \frac{\left(j(x,s)+D(\rho (x,s))\frac{d\rho (x,s)}{dx}\nu \sigma (\rho (x,s))\right)^2}{2\sigma (\rho (x,s))}}]`$
which is exactly (4,I). Furthermore (40) with the scaling (41,42) leads to (6).
Acknowledgments: We thank C. Bahadoran, L. Bertini, A. De Sole, D. Gabrielli, G. JonaโLasinio, C. Landim for very helpful discussions.
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# Adiabatic radiation reaction to the orbits in Kerr Spacetime
## I Introduction
A supermassive black hole (SMBH) accompanied by a compact object (CO) is among the most promising candidates for gravitational wave sources. This system may provide us the best opportunity for testing general relativity in the strong gravity regime. For this purpose, however, we need an accurate theoretical prediction of its waveform.
To investigate gravitational waves from SMBH-CO binary system, we use the black hole perturbation method. The background geometry is Kerr spacetime, and CO is described by a point particle. In the lowest order in mass ratio, the particle moves along the background geodesic. In the next order its orbit deviates from the geodesic due to radiation reaction effects.
In the Schwarzschild background, the particleโs orbit can be characterized solely by energy $`E`$ and azimuthal angular momentum $`L`$. We can evaluate the orbital evolution from the change rates of energy and azimuthal angular momentum, $`dE/dt`$ and $`dL/dt`$. In the adiabatic approximation the energy and the angular momentum that a particle loses are equated with the ones radiated to the infinity or into the black hole horizon as gravitational waves since there are conservation lows for $`E`$ and $`L`$ including the gravitational field. In this manner we can determine $`dE/dt`$ and $`dL/dt`$ from the asymptotic behavior of the radiated gravitational waves in the Schwarzschild case.
On the other hand, Carter constant $`Q`$ is also necessary to specify a geodesic in Kerr background. Since we do not have a conservation law corresponding to Carter constant, the change rate of $`Q`$ is not directly related to the asymptotic gravitational waves. Instead, we need to directly calculate the self-force acting on the particle Ori95 . Though the prescription to calculate the self-force is formally established MST97 ; QW97 , performing explicit calculation is not so straight forward.
Galโtsov Gal'tsov82 employed the *radiative* part of the metric perturbation, which was introduced earlier by Dirac Dirac38 , to calculate $`dE/dt`$ and $`dL/dt`$. He showed that this scheme correctly reproduces the standard results obtained by using the conservation laws. Recently, Mino gave a justification for applying the same scheme to $`dQ/dt`$ for bound orbits Mino03 . (See also Ref. Hughes05 .) However, actual implementation of $`dQ/dt`$ calculation again is not so straight forward. In this letter we derive a rather simple new formula for the adiabatic evolution of Carter constant.
## II Background
We consider the background Kerr spacetime in the Boyer-Lindquist coordinates: $`ds^2=(12Mr/\mathrm{\Sigma })dt^2(4Mar\mathrm{sin}^2\theta /\mathrm{\Sigma })dtd\phi +(\mathrm{\Sigma }/\mathrm{\Delta })dr^2+\mathrm{\Sigma }d\theta ^2+(r^2+a^2+2Ma^2r\mathrm{sin}^2\theta /\mathrm{\Sigma })\mathrm{sin}^2\theta d\phi ^2,`$ where $`\mathrm{\Sigma }=r^2+a^2\mathrm{cos}^2\theta `$ and $`\mathrm{\Delta }=r^22Mr+a^2.`$ Here $`M`$ and $`aM`$ are the mass and the angular momentum of the black hole, respectively. There are two Killing vectors $`\xi _{(t)}^\mu :=(_t)^\mu `$ and $`\xi _{(\phi )}^\mu :=(_\phi )^\mu `$ . In addition, Kerr spacetime possesses the Killing tensor, $`K_{\mu \nu }:=2\mathrm{\Sigma }l_{(\mu }n_{\nu )}+r^2g_{\mu \nu },`$ where the parentheses denote symmetrization on the indices enclosed, and $`l^\mu :=(r^2+a^2,\mathrm{\Delta },0,a)/\mathrm{\Delta }`$ and $`n^\mu :=(r^2+a^2,\mathrm{\Delta },0,a)/2\mathrm{\Sigma }`$ are two radial null vectors. Killing tensor satisfies the equation $`K_{(\mu \nu ;\rho )}=0`$,
We consider motion of a point particle, $`z^\alpha (\tau )=(t_z(\tau ),r_z(\tau ),\theta _z(\tau ),\varphi _z(\tau ))`$. Here $`\tau `$ is the proper time along the orbit. For geodesic motion, there are three constants of motion,
$`E`$ $`:=`$ $`u^\alpha \xi _\alpha ^{(t)}=\left(1{\displaystyle \frac{2Mr_z}{\mathrm{\Sigma }}}\right)u^t+{\displaystyle \frac{2Mar_z\mathrm{sin}^2\theta _z}{\mathrm{\Sigma }}}u^\phi ,`$ (1)
$`L`$ $`:=`$ $`u^\alpha \xi _\alpha ^{(\phi )}={\displaystyle \frac{2Mar_z\mathrm{sin}^2\theta _z}{\mathrm{\Sigma }}}u^t+{\displaystyle \frac{(r_z^2+a^2)^2\mathrm{\Delta }a^2\mathrm{sin}^2\theta _z}{\mathrm{\Sigma }}}\mathrm{sin}^2\theta _zu^\phi ,`$ (2)
$`Q`$ $`:=`$ $`K_{\alpha \beta }u^\alpha u^\beta ={\displaystyle \frac{(LaE\mathrm{sin}^2\theta _z)^2}{\mathrm{sin}^2\theta _z}}+a^2\mathrm{cos}^2\theta _z+\mathrm{\Sigma }^2(u^\theta )^2,`$ (3)
where $`u^\alpha :=dz^\alpha /d\tau `$. In addition, we define another notation for the Carter constant, $`C:=Q(aEL)^2`$. For orbits on the equatorial plane $`C`$ vanishes.
## III Geodesic Motion in Kerr Spacetime
Introducing a new parameter $`\lambda `$ defined by $`d\lambda =d\tau /\mathrm{\Sigma }`$, the geodesic equations become
$`\left({\displaystyle \frac{dr_z}{d\lambda }}\right)^2=R(r_z),\left({\displaystyle \frac{d\mathrm{cos}\theta _z}{d\lambda }}\right)^2=\mathrm{\Theta }(\mathrm{cos}\theta _z),`$ (4)
$`{\displaystyle \frac{dt_z}{d\lambda }}=a(aE\mathrm{sin}^2\theta _zL)+{\displaystyle \frac{r_z^2+a^2}{\mathrm{\Delta }}}P(r_z),{\displaystyle \frac{d\phi _z}{d\lambda }}=aE+{\displaystyle \frac{L}{\mathrm{sin}^2\theta _z}}+{\displaystyle \frac{a}{\mathrm{\Delta }}}P(r_z),`$ (5)
where $`P(r)=E(r^2+a^2)aL,R(r)=[P(r)]^2\mathrm{\Delta }[r^2+Q]`$ and $`\mathrm{\Theta }(\mathrm{cos}\theta )=C(C+a^2(1E^2)+L^2)\mathrm{cos}^2\theta +a^2(1E^2)\mathrm{cos}^4\theta .`$ It should be noted that the equation for the $`r`$-component and the one for the $`\theta `$-component are decoubled by using $`\lambda `$. Both $`R(r_z)`$ and $`\mathrm{\Theta }(\mathrm{cos}\theta _z)`$ are quartic functions of their arguments. Hence both solutions are given by elliptic functions. For bound orbits, we can systematically expand $`r_z`$ and $`\mathrm{cos}\theta _z`$ in Fourier series.
The other two equations (5) are integrated as
$`t_z(\lambda )`$ $`=`$ $`t^{(r)}(\lambda )+t^{(\theta )}(\lambda )+{\displaystyle \frac{dt_z}{d\lambda }}\lambda ,`$ (6)
$`\phi _z(\lambda )`$ $`=`$ $`\phi ^{(r)}(\lambda )+\phi ^{(\theta )}(\lambda )+{\displaystyle \frac{d\phi _z}{d\lambda }}\lambda ,`$ (7)
where $`\mathrm{}`$ means time average along the geodesic. $`t^{(r)}(\lambda ):=d\lambda \{(r_z^2+a^2)P(r_z)/\mathrm{\Delta }`$ $`(r_z^2+a^2)P(r_z)/\mathrm{\Delta }\}`$ and $`t^{(\theta )}(\lambda ):=d\lambda \{a(aE\mathrm{sin}^2\theta _zL)`$$`a(aE\mathrm{sin}^2\theta _zL)\},`$ are periodic functions with periods $`2\pi \mathrm{\Omega }_r`$ and $`2\pi \mathrm{\Omega }_\theta `$, respectively. Functions $`\phi ^{(r)}`$ and $`\phi ^{(\theta )}`$ are also defined in a similar way.
## IV Adiabatic evolution of constants of motion
In Ref. Mino03 , it was shown that the adiabatic radiation reaction to the constants of motion $`I^i=\{E,L,Q\}`$ can be evaluated by
$$\frac{dI^i}{d\lambda }=\underset{T\mathrm{}}{lim}\frac{1}{2T}_T^T๐\lambda \mathrm{\Sigma }\frac{I^i}{u^\alpha }f^\alpha [h_{\mu \nu }^{\mathrm{rad}}],$$
(8)
where $`h_{\mu \nu }^{\mathrm{rad}}`$ is the radiative part of the metric perturbation defined by half retarded field minus half advanced field, i.e., $`h_{\mu \nu }^{\mathrm{rad}}:=(h_{\mu \nu }^{\mathrm{ret}}h_{\mu \nu }^{\mathrm{adv}})/2.`$ Radiative field is a solution of source free vacuum Einstein equation. The singular parts contained in both retarded and advanced fields cancel out. Therefore we can avoid the tedious issue of regularizing self-force. $`f^\alpha `$ is a differential operator,
$$f^\alpha [h_{\mu \nu }]:=\frac{1}{2}(g^{\alpha \beta }+u^\alpha u^\beta )(h_{\beta \gamma ;\delta }+h_{\beta \delta ;\gamma }h_{\gamma \delta ;\beta })u^\gamma u^\delta .$$
### IV.1 Calculation of $`\dot{E}`$ and $`\dot{L}`$
It was shown by Galโtsov Gal'tsov82 that
$`{\displaystyle \frac{dE}{d\lambda }}`$ $`=`$ $`\underset{T\mathrm{}}{lim}{\displaystyle \frac{1}{2T}}{\displaystyle _T^T}๐\lambda \mathrm{\Sigma }(\xi _{(t)}^\alpha )f_\alpha [h_{\mu \nu }]`$ (9)
$`=`$ $`\underset{T\mathrm{}}{lim}{\displaystyle \frac{1}{2T}}{\displaystyle _T^T}๐\lambda \left[(\xi _{(t)}^\alpha )_\alpha \psi (x)\right]_{xz(\lambda )},`$ (10)
where $`\psi (x)=\mathrm{\Sigma }\stackrel{~}{u}^\mu \stackrel{~}{u}^\nu h_{\mu \nu }/2`$ and $`(\stackrel{~}{u}_t,\stackrel{~}{u}_r,\stackrel{~}{u}_\theta ,\stackrel{~}{u}_\phi ):=(E,\pm \sqrt{R(r)}/\mathrm{\Delta },\pm \sqrt{\mathrm{\Theta }(\mathrm{cos}\theta )}/\mathrm{sin}\theta ,L)`$. This vector field $`\stackrel{~}{u}_\mu `$ is an extension of the four velocity of a particle in the sense that it satisfies $`\stackrel{~}{u}_\mu (z(\lambda ))=u_\mu (\lambda )`$. Since in fact $`\stackrel{~}{u}_r`$ and $`\stackrel{~}{u}_\theta `$, respectively, depend only on $`r`$ and $`\theta `$, we can verify the relation, $`\stackrel{~}{u}_{\alpha ;\beta }=\stackrel{~}{u}_{\beta ;\alpha }`$.
Furthermore, Galโtsov has shown<sup>1</sup><sup>1</sup>1In Ref. Gal'tsov82 , not the in-field but the out-field was used. that $`\psi (x)`$ is given by
$`\psi (x)=i{\displaystyle \frac{d\omega }{2\pi \omega }\underset{\mathrm{},m}{}\varphi _{\omega ,\mathrm{},m}^{(in)}(x)๐\lambda ^{}\overline{\varphi _{\omega ,\mathrm{},m}^{(in)}(z(\lambda ^{}))}},`$ (11)
where $`\varphi _{\omega ,\mathrm{},m}(x)=\mathrm{\Sigma }\stackrel{~}{u}_\mu (x)\stackrel{~}{u}_\nu (x)\pi _{\omega ,\mathrm{},m}^{\mu \nu }(x)=\mathrm{\Sigma }\stackrel{~}{u}^\mu (r,\theta )\stackrel{~}{u}^\nu (r,\theta )\tau _{\mu \nu }^{}(r,\theta )e^{i\omega t+im\phi }\mathrm{\Delta }^2`$ $`{}_{2}{}^{}R_{\omega ,\mathrm{},m}^{}(r)_2S_{\omega ,\mathrm{},m}(\theta ).`$ $`\pi _{\omega ,\mathrm{},m}`$ is an appropriately normalized mode function of metric perturbations, which is constructed by applying a second rank differential operator $`\tau _{\mu \nu }^{}`$ to a mode function of Teukolsky equation Chrzan75 . The method to solve Teukolsky equation is well established Tagoshi . We can calculate the contribution from waves absorbed into a black hole by replacing the in-field to the up-field in Eq. (11).
We can express $`\varphi _{\omega ,\mathrm{},m}`$ as
$`\varphi _{\omega ,\mathrm{},m}(z(\lambda ))=e^{i(\omega t_z(\lambda )m\phi _z(\lambda ))}\mathrm{\Phi }_{\mathrm{},m}(r_z(\lambda ),dr_z/d\lambda (\lambda ),\theta _z(\lambda ),d\theta _z/d\lambda (\lambda )).`$ (12)
In the folloing text, we abbreviate $`dr_z/d\lambda `$ and $`d\theta _z/d\lambda `$ from the arguments for brevity. Here the exponent contains $`t^{(r)}(\lambda )`$, $`t^{(\theta )}(\lambda )`$, $`\phi ^{(r)}(\lambda )`$ and $`\phi ^{(\theta )}(\lambda )`$. Since $`r_z`$, $`t^{(r)}`$ and $`\phi ^{(r)}`$ ($`\theta _z`$, $`t^{(\theta )}`$ and $`\phi ^{(\theta )}`$) are periodic functions with period $`2\pi \mathrm{\Omega }_r^1`$ ($`2\pi \mathrm{\Omega }_\theta ^1`$), we can expand $`e^{i(\omega (t^{(r)}(\lambda _r)+t^{(\theta )}(\lambda _\theta ))m(\phi ^{(r)}(\lambda _r)+\phi ^{(\theta )}(\lambda _\theta )))}\mathrm{\Phi }_{\mathrm{},m}(r_z(\lambda _r),\theta _z(\lambda _\theta ))`$ into Fourier series as $`dt_z/d\lambda _{n_r,n_\theta }Z_{\omega ,\mathrm{},m}^{n_r,n_\theta }e^{in_r\mathrm{\Omega }_r\lambda _r+in_\theta \mathrm{\Omega }_\theta \lambda _\theta }`$ . Therefore we obtain $`\mathrm{exp}[i\omega (t^{(r)}(\lambda )+t^{(\theta )}(\lambda ))+im(\phi ^{(r)}(\lambda )+\phi ^{(\theta )}(\lambda ))]\mathrm{\Phi }_{\mathrm{},m}(r_z(\lambda ),\theta _z(\lambda ))=dt_z/d\lambda _{n_r,n_\theta }Z_{\omega ,\mathrm{},m}^{n_r,n_\theta }e^{i(n_r\mathrm{\Omega }_r+n_\theta \mathrm{\Omega }_\theta )\lambda }.`$ Using this expansion, we obtain
$`{\displaystyle ๐\lambda ^{}\overline{\varphi _{\omega ,\mathrm{},m}^{(in)}(z(\lambda ^{}))}}`$ $`=`$ $`{\displaystyle \underset{n_r,n_\theta }{}}2\pi \delta \left(\omega \omega _m^{n_r,n_\theta }\right)\overline{Z_{\mathrm{},m}^{n_r,n_\theta }},`$ (13)
where $`\omega _m^{n_r,n_\theta }=dt_z/d\lambda ^1\left(md\phi _z/d\lambda +n_r\mathrm{\Omega }_r+n_\theta \mathrm{\Omega }_\theta \right)`$ and $`Z_{\mathrm{},m}^{n_r,n_\theta }:=Z_{\omega _m^{n_r,n_\theta },\mathrm{},m}^{n_r,n_\theta }`$. Substituting Eq. (13) into Eq. (10) with Eq. (11), and integrating it over $`\omega `$, we obtain
$`{\displaystyle \frac{dE}{d\lambda }}={\displaystyle \frac{dt_z}{d\lambda }}\underset{T\mathrm{}}{lim}{\displaystyle \frac{1}{2T}}{\displaystyle _T^T}๐\lambda {\displaystyle \underset{\mathrm{},m,}{}_{n_r,n_\theta ,n_r^{},n_\theta ^{}}}Z_{\mathrm{},m}^{n_r,n_\theta }\overline{Z_{\mathrm{},m}^{n_r^{},n_\theta ^{}}}e^{i((n_rn_r^{})\mathrm{\Omega }_r+(n_\theta n_\theta ^{})\mathrm{\Omega }_\theta )\lambda }.`$ (14)
Now integration over $`\lambda `$ is straight forward. We finally end up with the well known formula except for the overall normalization depending on the definition of the mode function:
$$\frac{dE}{dt}=\underset{\mathrm{},m,n_r,n_\theta }{}|Z_{\mathrm{},m}^{n_r,n_\theta }|^2.$$
(15)
In a similar manner, a formula for the angular momentum loss rate is obtained as
$$\frac{dL}{dt}=\underset{\mathrm{},m,n_r,n_\theta }{}\frac{m}{\omega _m^{n_r,n_\theta }}|Z_{\mathrm{},m}^{n_r,n_\theta }|^2.$$
(16)
### IV.2 calculation of $`\dot{Q}`$
The expression for the radiation reaction to the Carter constant can be cast into a form analogous to the energy loss rate and the angular momentum loss rate. Substituting
$$f_\nu =\frac{1}{2}(_\nu h_{\alpha \beta })u^\alpha u^\beta \frac{d}{d\tau }(h_{\nu \beta }u^\beta )\frac{1}{2}u_\nu \frac{d}{d\tau }(u^\alpha u^\beta h_{\alpha \beta }),$$
(17)
the evolution of Carter constant is given by
$`{\displaystyle \frac{dQ}{d\tau }}=2K_\mu ^\nu u^\mu f_\nu 2\left[K_\mu ^\nu \stackrel{~}{u}^\mu _\nu {\displaystyle \frac{\psi }{\mathrm{\Sigma }}}+h_{\alpha \beta }\stackrel{~}{u}^\alpha \stackrel{~}{u}^\mu (K_{\mu ;\nu }^\beta \stackrel{~}{u}^\nu K_\mu ^\nu \stackrel{~}{u}_{;\nu }^\beta )\right]_{xz(\lambda )}.`$ (18)
Here โ$``$โ means that terms which become a total derivative or $`O(h^2)`$ are neglected. Furthermore we can show that the second term in the above equaiton vanishes by using the facts $`\stackrel{~}{u}_{\alpha ;\beta }=\stackrel{~}{u}_{\beta ;\alpha }`$ and $`K_{(\mu \nu ;\rho )}=0`$. Thus we obtain
$$\frac{dQ}{d\lambda }=\underset{T\mathrm{}}{lim}\frac{1}{T}_T^T๐\lambda \mathrm{\Sigma }K_\mu ^\nu \stackrel{~}{u}^\mu _\nu \frac{\psi (x)}{\mathrm{\Sigma }},$$
(19)
which is written in terms of $`\psi (x)`$. Hence, we find that the change rate of Carter constant is obtained by replacing $`\xi _{(t)}^\alpha _\alpha `$ in the above expression for $`dE/dt`$ given in Eq. (11) with $`2\mathrm{\Sigma }K_\mu ^\nu \stackrel{~}{u}^\mu _\nu \mathrm{\Sigma }^1.`$ Then, what we have to evaluate is
$`{\displaystyle }`$ $`d\lambda `$ $`\left[\mathrm{\Sigma }K_\mu ^\nu \stackrel{~}{u}^\mu _\nu {\displaystyle \frac{\psi (x)}{\mathrm{\Sigma }}}\right]_{xz(\lambda )}`$ (20)
$`=`$ $`{\displaystyle ๐\lambda \left[\mathrm{\Sigma }\left(\mathrm{\Sigma }(\mathrm{}^\mu \stackrel{~}{u}_\mu n^\nu _\nu +n^\nu \stackrel{~}{u}_\mu \mathrm{}^\mu _\nu )+r^2\stackrel{~}{u}^\mu _\mu \right)\frac{\psi (x)}{\mathrm{\Sigma }}\right]_{xz(\lambda )}}`$ (21)
$`=`$ $`{\displaystyle ๐\lambda \left[\left(\frac{P(r)}{\mathrm{\Delta }}((r^2+a^2)_t+a_\phi )\frac{dr_z}{d\lambda }_r\right)\psi (x)\right]_{xz(\lambda )}}.`$ (22)
In the last step the last term have been integrated by parts using $`\stackrel{~}{u}^\mu _\mu =\mathrm{\Sigma }^1d/d\lambda `$, which is valid after substitution of $`z(\lambda )`$.
## V Further reduction
Further simplification is possible. For an arbitrary function of $`r_z`$ and $`\theta _z`$, we have
$`\underset{T\mathrm{}}{lim}{\displaystyle \frac{1}{2T}}{\displaystyle _T^T}๐\lambda \mathrm{exp}\left[i\omega _m^{n_r,n_\theta }t_z(\lambda )+im\phi _z(\lambda )\right]f(r_z(\lambda ),\theta _z(\lambda ))`$ (23)
$`={\displaystyle \frac{\mathrm{\Omega }_r\mathrm{\Omega }_\theta }{(2\pi )^2}}{\displaystyle _0^{2\pi \mathrm{\Omega }_r^1}}๐\lambda _r{\displaystyle _0^{2\pi \mathrm{\Omega }_\theta ^1}}๐\lambda _\theta \mathrm{exp}\left[in_r\mathrm{\Omega }_r\lambda _ri\omega _m^{n_r,n_\theta }t^{(r)}(\lambda _r)+im\phi ^{(r)}(\lambda _r)\right]`$ (24)
$`\times \mathrm{exp}\left[in_\theta \mathrm{\Omega }_\theta \lambda _\theta i\omega _m^{n_r,n_\theta }t^{(\theta )}(\lambda _\theta )+im\phi ^{(\theta )}(\lambda _\theta )\right]f(r_z(\lambda _r),\theta _z(\lambda _\theta )).`$ (25)
This relation can be easily verified by substituting $`\mathrm{exp}[i\omega _m^{n_r,n_\theta }(t^{(r)}(\lambda _r)+t^{(\theta )}(\lambda _\theta ))+im(\phi ^{(r)}(\lambda _r)+\phi ^{(\theta )}(\lambda _\theta )))]f(r_z(\lambda _r),\theta _z(\lambda _\theta ))=_{n_r,n_\theta }f_{n_r,n_\theta }`$ $`e^{i(n_r\mathrm{\Omega }_r+n_\theta \mathrm{\Omega }_\theta )\lambda }`$. Then, by using the above relation, the $`\lambda `$-integral in Eq. (22) can be decomposed. The part of $`\lambda _r`$-integral takes the following form, and it can be integrated by parts as
$`{\displaystyle ๐\lambda _r\left[\frac{dr_z}{d\lambda }_r\mathrm{exp}\left[in_r\mathrm{\Omega }_r\lambda _ri\omega t+im\phi \right]f(r,\theta )\right]_{rr_z(\lambda _r),tt^{(r)}(\lambda _r),\phi \phi ^{(r)}(\lambda _r)}}`$ (26)
$`={\displaystyle }d\lambda _r\{[{\displaystyle \frac{d}{d\lambda _r}}\left({\displaystyle \frac{dt^{(r)}}{d\lambda _r}}\right)_t\left({\displaystyle \frac{d\phi ^{(r)}}{d\lambda _r}}\right)_\phi +in_r\mathrm{\Omega }_r]`$ (27)
$`\times \mathrm{exp}[in_r\mathrm{\Omega }_r\lambda _ri\omega t+im\phi ]f(r,\theta )\}_{rr_z(\lambda _r),tt^{(r)}(\lambda _r),\phi \phi ^{(r)}(\lambda _r)}.`$ (28)
The first term in the last line is a total derivative, and therefore does not contribute to the average over a long period of time. Combining the above results, we obtain
$`{\displaystyle \frac{dQ}{dt}}`$ $`=`$ $`{\displaystyle \frac{dt_z}{d\lambda }}^1\underset{T\mathrm{}}{lim}{\displaystyle \frac{1}{T}}{\displaystyle _T^T}\left(\left[{\displaystyle \frac{(r^2+a^2)P}{\mathrm{\Delta }}}_t+{\displaystyle \frac{aP}{\mathrm{\Delta }}}_\phi +in_r\mathrm{\Omega }_r\right]\psi (x)\right)_{xz(\lambda )}๐\lambda `$ (29)
$`=`$ $`2{\displaystyle \frac{(r^2+a^2)P}{\mathrm{\Delta }}}{\displaystyle \frac{dE}{dt}}2{\displaystyle \frac{aP}{\mathrm{\Delta }}}{\displaystyle \frac{dL}{dt}}+2{\displaystyle \underset{\mathrm{},m,n_r,n_\theta }{}}{\displaystyle \frac{n_r\mathrm{\Omega }_r}{\omega _m^{n_r,n_\theta }}}|Z_{\mathrm{},m}^{n_r,n_\theta }|^2.`$ (30)
This expression is as easy to evaluate as $`dE/dt`$ and $`dL/dt`$. To evaluate the last term, we have only to replace $`m`$ in the expression for $`dL/dt`$ with $`n_r\mathrm{\Omega }_r`$.
## VI Consistency check
We know that a circular orbit stays circular under radiation reaction Kennefick . This condition becomes $`dQ/dt=(2(r^2+a^2)P/\mathrm{\Delta })dE/dt(2aP/\mathrm{\Delta })`$ $`dL/dt.`$ For circular orbits $`Z_{\mathrm{},m}^{n_r,n_\theta }=0`$ when $`n_r0`$. Therefore the last term in Eq. (30) vanishes. Then Eq. (30) agrees with the above condtion that a circular orbit stays circular.
We also know that an orbit on the equatorial plane does not leave the equatorial plane by symmetry. This can be clearly seen by rewriting the above formula in terms of $`C`$. An identity $`dt_z/d\lambda dE/dtd\phi _z/d\lambda dL/dt+_{\mathrm{},m,n_r,n_\theta }\{(n_r\mathrm{\Omega }_r+n_\theta \mathrm{\Omega }_\theta )/\omega _m^{n_r,n_\theta }\}|Z_{\mathrm{},m}^{n_r,n_\theta }|^2=0`$ follows from the definition of $`\omega _m^{n_r,n_\theta }`$ with the aid of the expressions for $`dE/dt`$ and $`dL/dt`$ given in Eqs. (15) and (16). Using this identity, we have
$`{\displaystyle \frac{dC}{dt}}`$ $`=`$ $`2a^2E\mathrm{cos}^2\theta _z{\displaystyle \frac{dE}{dt}}2L\mathrm{cot}^2\theta _z{\displaystyle \frac{dL}{dt}}2{\displaystyle \underset{\mathrm{},m,n_r,n_\theta }{}}{\displaystyle \frac{n_\theta \mathrm{\Omega }_\theta }{\omega _m^{n_r,n_\theta }}}|Z_{\mathrm{},m}^{n_r,n_\theta }|^2.`$ (31)
From this equation it is manifest that $`dC/dt=0`$ when $`\theta =\pi /2`$. Notice that $`Z_{\mathrm{},m}^{n_r,n_\theta }0`$ only for $`n_\theta =0`$ in the case of equatorial orbits.
###### Acknowledgements.
We would like to thank Y. Mino for useful discussions. This work is supported in part by JSPS Research Fellowships for Young Scientists, Nos. 5919 and 1756, by Grant-in-Aid for Scientific Research, No. 16740141 from Japan Society for the Promotion of Scient (JSPS), by Grant-in-Aid for Scientific Research, Nos. 14047212, 14047214, and by that for the 21st Century COE by that for the 21st Century COE โCenter for Diversity and Universality in Physicsโ at Kyoto university from the Ministry of Education, Culture, Sports, Science and Technology (MEXT) of Japan.
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# 1 Introduction
## 1 Introduction
During the last decades the tests of the perturbative QCD are focused on the experimental \[1-15\] and theoretical study \[16-26\] of hard processes like direct photon and $`\pi ^0`$-production with large transverse momentum in $`pp`$, $`pA`$ and $`AA`$ collisions. These processes offer a unique possibility to determine the parton distribution functions (PDF) inside the proton as well as the parton fragmentation functions (PFF). In particular the good understanding of the $`pp`$ data is the prerequisite for any attempt to extract new physics, related to the formation of a quark-gluon plasma phase, from the $`pA`$ and $`AA`$ data. Extensive studies of the $`\pi `$ (or $`\gamma `$) production in $`pp`$ collisions have shown that the transverse momentum distribution $`g(k_T)`$ of the partons inside the proton has to be taken into account for a succesfull description of the observed $`p_T`$ spectrum \[20-23\]. In all these investigations one assumes a Gaussian form for $`g(k_T)`$. A new nonperturbative parameter is introduced through this approach: the mean instrinsic transverse momentum $`<k_T>`$ of the partons. Although the data of some experiments concerning hadron production at large $`p_T`$ could be explained with relative small $`<k_T>`$ values ($`0.30.5GeV`$), compatible with the Heisenberg uncertainty relation for partons inside the proton, there are a number of other processes leading to a large mean transverse momentum ($`<k_T>14GeV`$), depending on $`Q^2`$, for a description of the corresponding experimental data. Such a value of $`<k_T>`$ is too high and cannot be explained as an internal structure of the proton . However, as mentioned by several authors the form of $`g(k_T)`$ can influence significantly the value of $`<k_T>`$ as well as its $`p_T`$ dependence. In the present work we derive a transverse momentum distribution for the partons inside the proton using a potential quark model which has succesfully been used to describe the spectra of mesonic and baryonic bound states in the past. Following we investigate the three-body quantum mechanical bound state problem solving numerically the Schrรถdinger equation and obtaining the single particle transverse momentum distribution for the consistuent parton. Our main assumption is that intrinsic transverse momentum effects are not influenced by the Lorenz boost along the beam axis and therefore could be treated within a nonrelativistic approach. We use the derived distribution to fit experimental data for the $`pp\pi _0+X`$ process. In particular we investigate the measurements for the $`p_T`$-spectrum of the outcoming $`\pi ^0`$ in three experiments performed at different center of mass energies. It turns out that a relatively low mean transverse momentum $`<k_T>`$ ($`O(300MeV)`$), compatible with intrinsic dynamics inside the proton, for the initial partons is sufficient in order to fit perfectly the experimental data. A smooth dependence of $`<k_T>`$ on $`p_T`$, which within our approach is induced by a corresponding variation of the string tension in the quark-quark potential, is required. Our analysis shows that the $`k_T`$ distribution of the consistuent partons can be strongly influenced by three-body effects and the form of the confining potential which have to be taken into account in order to describe correctly the experimental data concerning the pion production in $`pp`$ collisions. The paper is organized as follows: in Section 2 we present the parton model differential cross section as well as the corresponding kinematics for the $`\pi ^0`$-production in $`pp`$-collisions. In section 3 we derive the intrinsic transverse momentum disrtibution for the partons inside the proton using the quark potential model of . In section 4 we present our numerical results concerning the description of the data of three different experiments as well as the corresponding dependence $`<k_T(p_T)>`$. Finally in section 5 we summarize our study and we discuss possible extensions of the present analysis.
## 2 The $`pp\pi ^0+X`$ cross section
In the lowest-order perturbative QCD (pQCD), the differential cross section for the direct pion production in $`pp`$ collisions is given by:
$`E_\pi {\displaystyle \frac{d\sigma }{d^3p}}(pp\pi ^0+X)`$ $`=`$ $`K{\displaystyle \underset{abcd}{}}{\displaystyle ๐x_a๐x_bf_{a/p}(x_a,Q^2)f_{b/p}(x_b,Q^2)}`$ (1)
$`\times `$ $`{\displaystyle \frac{d\sigma }{d\widehat{t}}}(abcd){\displaystyle \frac{D_{\pi /c}(z_c,Q^2)}{\pi z_c}}`$
where $`f_{i/p}`$ ($`i=a,b`$) are the longitudinal parton distribution functions (PDF) for the colliding partons $`a`$ and $`b`$ while $`D_{\pi /c}`$ is the parton fragmentation function (PFF) for the pion. For the scale $`Q`$ we use the relation $`Q^2=\frac{2\widehat{s}\widehat{t}\widehat{u}}{\widehat{s}^2+\widehat{t}^2+\widehat{u}^2}`$ proposed in with $`\widehat{s},\widehat{t},\widehat{u}`$ the usual Mandelstam variables. The variable $`z_c`$ indicates the momentum fraction carried by the final hadron. The higher order corrections are taken into account by choosing $`K2`$ for the corresponding coefficient in (1) in the $`p_T`$ region of interest. It is straightforward to include partonic transverse degrees of freedom using the following replacement in the PDFs of eq.(1):
$$dx_if_{i/p}(x_i,Q^2)dx_id^2k_{T,i}g(\stackrel{}{k}_{T,i})f_{i/p}(x_i,Q^2)$$
(2)
with $`i=a,b`$. In order to avoid singularities in the differential cross sections describing the partonic subprocesses we introduce a regularizing parton mass $`m=0.8GeV`$, as in , in the Mandelstam variables occuring in the denominator of the coresponding matrix elements. The explicit formulas of the relevant partonic cross sections can be found in . Taking into account the transverse degrees of freedom we get the following expressions for the variables $`\widehat{s},\widehat{t},\widehat{u}`$:
$`\widehat{s}`$ $`=`$ $`sx_ax_b+{\displaystyle \frac{k_{T,a}^2k_{T,b}^2}{sx_ax_b}}2\stackrel{}{k}_{T,a}\stackrel{}{k}_{T,b}`$
$`\widehat{t}`$ $`=`$ $`(x_a+{\displaystyle \frac{k_{T,a}^2}{sx_a}}){\displaystyle \frac{p_T\sqrt{s}}{z_c}}+{\displaystyle \frac{2}{z_c}}\stackrel{}{k}_{T,a}\stackrel{}{p}_T`$
$`\widehat{u}`$ $`=`$ $`(x_b+{\displaystyle \frac{k_{T,b}^2}{sx_b}}){\displaystyle \frac{p_T\sqrt{s}}{z_c}}+{\displaystyle \frac{2}{z_c}}\stackrel{}{k}_{T,b}\stackrel{}{p}_T`$ (3)
Due to energy-momentum conservation the momentum fraction of the final hadron $`z_c`$ is given by:
$$z_c=\frac{(x_a+\frac{k_{T,a}^2}{sx_a}+x_b+\frac{k_{T,b}^2}{sx_b})p_T\sqrt{s}2(\stackrel{}{k}_{T,a}+\stackrel{}{k}_{T,b})\stackrel{}{p}_T}{\widehat{s}}$$
(4)
For a consistent description of the kinematics in the partonic subprocesses we imply the cuts:
$$z_c1;k_{T,a}^2<p_T\sqrt{s};k_{T,b}^2<p_T\sqrt{s}$$
(5)
To calculate the cross section given in eq.(1) we have first to determine the distribution $`g(\stackrel{}{k}_T)`$ and then perform the corresponding phase space integrations. Contrary to the usual treatment assuming a Gaussian form for $`g(\stackrel{}{k}_T)`$ we will here derive an alternative expression based on a widely applied quark potential model.
## 3 The intrinsic transverse momentum distribution $`g(\stackrel{}{k}_T)`$
Since the early days of quantum chromodynamics the main and almost unique tool used for the description of the hadronic bound states (mesons, baryons) remain potential models for the quark-quark and quark-antiquark pair interaction. One of the most succesfull models was proposed by A. Martin used first to describe mesonic states and later to describe baryons . The quark-quark interaction within this model is given as:
$$V(r)=A_{qq}r^{0.1}+B_{qq}$$
(6)
and a similar expression with adapted coefficients $`A_{\overline{q}q}`$, $`B_{\overline{q}q}`$ holds for the quark-antiquark interaction. In the following we will consider a baryon consisting of 3 valence quarks interacting pairwise with the potential (6). The Hamiltonian operator of the system is given as:
$$\widehat{H}=\frac{\mathrm{}^2}{2m}(_1^2+_2^2+_3^2)+V(\stackrel{}{r}_{12})+V(\stackrel{}{r}_{23})+V(\stackrel{}{r}_{31})$$
(7)
The baryonic ground state can then be obtained by solving the corresponding Schrรถdinger equation. Following it is convenient to introduce Jacobi coordinates $`\stackrel{}{\xi }_i`$:
$`\stackrel{}{\xi }_1`$ $`=`$ $`\stackrel{}{r}_2\stackrel{}{r}_1`$
$`\stackrel{}{\xi }_2`$ $`=`$ $`{\displaystyle \frac{2\stackrel{}{r}_3\stackrel{}{r}_2\stackrel{}{r}_1}{\sqrt{3}}}`$
$`\stackrel{}{\xi }_3`$ $`=`$ $`{\displaystyle \frac{\stackrel{}{r}_1+\stackrel{}{r}_2+\stackrel{}{r}_3}{3}}`$ (8)
getting the equation:
$$\left[\frac{\mathrm{}^2}{2M}\stackrel{}{}_{\xi _3}^2\frac{\mathrm{}^2}{m}(\stackrel{}{}_{\xi _1}^2+\stackrel{}{}_{\xi _2}^2)+\stackrel{~}{V}(\stackrel{}{\xi }_1,\stackrel{}{\xi }_2)\right]\mathrm{\Psi }_G(\stackrel{}{\xi }_1,\stackrel{}{\xi }_2,\stackrel{}{\xi }_3)=E_G\mathrm{\Psi }_G(\stackrel{}{\xi }_1,\stackrel{}{\xi }_2,\stackrel{}{\xi }_3)$$
(9)
with $`M=3m`$ and $`m`$ is the constituent quark mass. Eliminating the translational mode ($`\stackrel{}{\xi }_3`$) we obtain a partial differential equation (PDE) depending solely on the variables $`\stackrel{}{\xi }_1`$, $`\stackrel{}{\xi }_2`$. The potential energy $`\stackrel{~}{V}`$ in eq.(9) is given as:
$$\stackrel{~}{V}(\stackrel{}{\xi }_1,\stackrel{}{\xi }_2)=V(\stackrel{}{\xi }_1)+V(\frac{1}{2}(\stackrel{}{\xi }_1\sqrt{3}\stackrel{}{\xi }_2))+V(\frac{1}{2}(\sqrt{3}\stackrel{}{\xi }_2+\stackrel{}{\xi }_1))$$
(10)
Introducing the radial variable $`\xi =\sqrt{\xi _1^2+\xi _2^2}`$ and the angular variables $`\chi ,\theta _1,\varphi _1,\theta _2,\varphi _2`$ we rewrite eq.(9) as follows:
$`{\displaystyle \frac{\mathrm{}^2}{m}}\{{\displaystyle \frac{1}{\xi ^5}}_\xi \left(\xi ^5_\xi \stackrel{~}{\mathrm{\Psi }}\right)+{\displaystyle \frac{1}{\xi ^2}}\left[{\displaystyle \frac{1}{\mathrm{sin}^22\chi }}_\chi \left(\mathrm{sin}^22\chi _\chi \stackrel{~}{\mathrm{\Psi }}\right)+{\displaystyle \frac{\widehat{L}_1^2\stackrel{~}{\mathrm{\Psi }}}{\mathrm{cos}^2\chi }}+{\displaystyle \frac{\widehat{L}_2^2\stackrel{~}{\mathrm{\Psi }}}{\mathrm{sin}^2\chi }}\right]\}`$
$`+\stackrel{~}{V}(\xi ,\chi ,\theta _1,\varphi _1,\theta _2,\varphi _2)\stackrel{~}{\mathrm{\Psi }}=E_G\stackrel{~}{\mathrm{\Psi }}`$ (11)
where the angular momentum operators are:
$$\widehat{L}_i^2=\left[\frac{1}{\mathrm{sin}\theta _i}\frac{}{\theta _i}(\mathrm{sin}\theta _i\frac{}{\theta _i})+\frac{1}{\mathrm{sin}^2\theta _i}\frac{^2}{\varphi _i^2}\right]$$
The reduced ground state wave function $`\stackrel{~}{\mathrm{\Psi }}`$ can be expressed in terms of the hyperspherical harmonics $`P_L(\mathrm{\Omega })`$ as follows:
$$\stackrel{~}{\mathrm{\Psi }}(\xi ,\chi ,\theta _1,\varphi _1,\theta _2,\varphi _2)=\underset{L=0}{\overset{\mathrm{}}{}}\frac{u_L(\xi )}{\xi ^{5/2}}P_L(\chi ,\theta _1,\varphi _1,\theta _2,\varphi _2)$$
(12)
The radial part of the wave function $`\stackrel{~}{\mathrm{\Psi }}`$ fullfils the equation :
$$\frac{d^2u_L}{d\xi ^2}\frac{15/4+L(L+4)}{\xi ^2}u_L+\frac{m}{\mathrm{}^2}E_Gu_L\frac{m}{\mathrm{}^2}\underset{L^{}}{}u_L^{}\stackrel{~}{V}_{LL^{}}=0$$
(13)
As mentioned in the matrix $`\stackrel{~}{V}_{LL^{}}=๐\mathrm{\Omega }P_L^{}(\mathrm{\Omega })\stackrel{~}{V}(\xi ,\mathrm{\Omega })P_L^{}(\mathrm{\Omega })`$ for the particular choise of the potential (6) is diagonal dominant and therefore to a very good approximation the ground state of the baryonic system is determined by the $`L=0`$ term in the expansion (12). Therefore the radial part of the ground state wavefunction obeys, within the above approximation, the equation:
$$\frac{d^2u_0}{d\xi ^2}\frac{15}{4\xi ^2}u_0+\frac{m}{\mathrm{}^2}(E_G\stackrel{~}{V}_{00})u_0=0$$
(14)
where: $`\stackrel{~}{V}_{00}=A_{00}+B_{00}\xi ^{0.1}`$ with $`A_{00}=\frac{3}{2}A_{qq}`$ and $`B_{00}=\frac{1}{2}\lambda B_{qq}`$. The constant $`\lambda `$ is given by: $`\lambda =\frac{24}{\pi }\mathrm{\Gamma }(1.55)\mathrm{\Gamma }(1.5)\mathrm{\Gamma }(3.05)`$. It is straightforward to show that the ground state wavefunction in the momentum space (conjugate to the space of $`\stackrel{}{\xi }_i,i=1,2,3`$) is given by:
$$\stackrel{~}{\mathrm{\Phi }}(k_\xi ^2)=N_0^{\mathrm{}}๐\xi \xi ^{1/2}u_0(\xi )\frac{J_2(k_\xi \xi )k_\xi \xi J_3(k_\xi \xi )}{k_\xi ^2}$$
(15)
where $`J_n`$ is the Bessel function of order $`n`$ and $`N`$ is a normalization constant. It is useful to determine the transformation of the momenta $`\stackrel{}{k}_{\xi _i}`$ to the cartesian momenta $`\stackrel{}{k}_i`$:
$`\stackrel{}{k}_{\xi _1}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\stackrel{}{k}_1\stackrel{}{k}_2)`$
$`\stackrel{}{k}_{\xi _2}`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{3}}}(\stackrel{}{k}_1+\stackrel{}{k}_22\stackrel{}{k}_3)`$
$`\stackrel{}{k}_{\xi _3}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{3}}}(\stackrel{}{k}_1+\stackrel{}{k}_2+\stackrel{}{k}_3)`$ (16)
The above expressions simplify in the center of mass frame $`\stackrel{}{k}_{\xi _3}=\stackrel{}{0}`$ of the baryonic system where the following relation is valid:
$$k_\xi ^2=k_1^2+k_2^2+\stackrel{}{k}_1\stackrel{}{k}_2$$
The two-particle density $`\rho (\stackrel{}{k}_1,\stackrel{}{k}_2)`$ in this case is given by:
$$\rho (\stackrel{}{k}_1,\stackrel{}{k}_2)=|\stackrel{~}{\mathrm{\Phi }}(k_\xi ^2)|^2=|\stackrel{~}{\mathrm{\Phi }}(k_1^2+k_2^2+\stackrel{}{k}_1\stackrel{}{k}_2)|^2$$
(17)
From eq.(17) we obtain the one-particle transverse momentum density $`g(\stackrel{}{k}_T)`$ as:
$$g(\stackrel{}{k}_T)=4\pi _0^{\mathrm{}}๐k_z_1^1๐z_0^{\mathrm{}}๐k_2k_2^2|\stackrel{~}{\mathrm{\Phi }}(k_T^2+k_z^2+k_2^2+zk_2\sqrt{k_T^2+k_z^2})|^2$$
(18)
where $`z`$ is the cosine of the angle between the vectors $`\stackrel{}{k}=(\stackrel{}{k}_T,k_z)`$ and $`\stackrel{}{k}_2`$. In fact the function $`u_0(\xi )`$ can be obtained only numerically by solving equation (14) using the Numerov algorithm. Therefore also the transverse momentum distribution (18) is known only numerically. The integrations in eqs.(15,18) can be performed to a great accuracy (relative error $`10^6`$) using a mixture of Gauss-Kronrod quadrature and the VEGAS Monte-Carlo integration routine . In Fig. 1 we present the density in transverse momentum space of a parton inside the proton obtained by our approach. The values of the constants $`A_{qq}`$ and $`B_{qq}`$ are choosen according to in order to fit the size and the binding energy of the proton. The characteristic structure dominating the form of $`g(\stackrel{}{k}_T)`$ within our model is the second maximum at relative high transverse momenta. This leads, as we will see in the next section, to a reduction, relative to the Gaussian case, of the mean intrinsic transverse momentum of the partons needed to describe the experimental data for the $`\pi ^0`$-production at various energies.
## 4 Numerical results
In the following we will use the derived distribution $`g(\stackrel{}{k}_T)`$ in order to calculate within the parton model the differential cross section for the $`\pi ^0`$-production in $`pp`$ collisions according to eq.(1). For the longitudinal parton distribution functions we use the recent Martin, Roberts, Stirling and Thorne (MRST) scheme while for the parton fragmentation functions we use the Kniehl, Kramer, Potter parametrization. The additional nonperturbative parameter in our treatment is the mean value of the intrinsic transverse momentum $`<k_T>`$ which is related to the string tension $`B_{qq}`$ in eq.(6). Although the distribution $`g(\stackrel{}{k}_T)`$ is derived for the valence quarks inside the proton here we will use the same distribution also for the initial gluons assuming universality at the level of consistuent partons. In any case for longitudinal momentum fraction $`x_i>0.5`$ ($`i=a,b`$) the contribution of valence quarks is dominant and our description is accurate. The phase space integrations are performed using the VEGAS Monte-Carlo routine. In order to fit the experimental data we allow $`<k_T>`$ to vary as a function of the beam energy and the transverse momentum ($`p_T`$) of the finally produced hadron. We will analyse here the results of three experiments concerning $`\pi ^0`$-production at different energies. The first set of data are taken from the fixed target experiment performed in Fermilab (protons incident on $`H_2`$ target) . The cross section for the $`\pi ^0`$-production with transverse momentum $`p_T`$ at midrapidity and for three different proton beam energies $`E_L=200,300`$ and $`400GeV`$ is measured. In Fig. 2 we present the various datasets using symbols while with solid lines we show the results of our calculation and with dashed lines the corresponding results using a Gaussian $`g(\stackrel{}{k}_T)`$. As we can see the two descriptions differ only in the region $`p_T<1GeV`$. Within our approach we need an intrinsic transverse momentum $`<k_T>`$ of the order of at most $`0.5GeV`$ in order to reproduce all the experimental data. It must be noted that in this case one can perfectly fit the data also for $`p_T<1GeV`$ which, as we see in Fig. 2, is not possible using a Gaussian distribution $`g(\stackrel{}{k}_T)`$. In Fig. 3 we show the functions $`<k_T>(p_T)`$ obtained using the quark model inspired function $`g(\stackrel{}{k}_T)`$ as well as a Gaussian form. It is evident that the Gaussian model leads to much higher values of $`<k_T>`$. In particular this difference is large even for large values of $`p_T`$ where our approach is more precise as the valence quark contribution is dominant. As already discussed in the previous section this difference relies on the fact that the non-Gaussian $`g(\stackrel{}{k}_T)`$ derived here posseses a second small local maximum at high transverse momenta (see Fig. 1) attributed to the form of the quark-quark potential and the many-body character of the system.
The second experiment we consider is the WA70 at CERN SPS . It is also a fixed target experiment with $`E_L=280GeV`$. We are interested in $`\pi ^0`$-production. In Fig. 4 we show the experimental data (full stars) and the corresponding pQCD calculation using the quark model inspired $`g(\stackrel{}{k}_T)`$ (solid line) as well as a Gaussian form (dashed line). Both distributions reproduce perfectly the experimental data and cannot be distinguished graphically. However our approach leads to significantly lower values of $`<k_T>`$ than in the Gaussian case. This can be clearly seen in Fig. 5 where we show the function $`<k_T(p_T)>`$ both for the quark potential model inspired $`g(\stackrel{}{k}_T)`$ (full circles) as well as the Gaussian intrinsic transverse momentum distribution (full stars). Also in this case the difference in $`<k_T>`$ between the two approaches is large in a $`p_T`$-region where the main contribution is attributed to the valence quarks. For comparison we present also the same function for the Fermilab experiment at $`E_L=300Gev`$ (open circles, see Fig. 3).
Finally we have analysed the $`pp\pi ^0+X`$ data of the most recent PHENIX experiment at the Relativistic Heavy Ion Collider with $`\sqrt{s}=200GeV`$ . The corresponding cross section can be described to a good accuracy without inclusion of any $`k_T`$-smearing, a fact which is compatible with expectations for the perturbative character of the subprocesses involved in this case. Here we have fitted the experimental data using non-Gaussian $`k_T`$-smearing effects. In this way we get a perfect description of the measured cross section. The obtained mean intrinsic transverse momentum is almost constant: $`<k_T>250MeV`$. Based on Heisenberg uncertainty relation one could explain this value of the mean transverse momentum as an effect of the internal partonic structure of the incident proton. Our results are presented in Figs. 6,7. The PHENIX data (open cirlces) for $`E\frac{d^3\sigma }{dp^3}`$ together with the parton model calculations (crosses) using non-Gaussian $`g(\stackrel{}{k}_T)`$ are shown in Fig. 6. The corresponding function $`<k_T>(p_T)`$ is presented in Fig. 7.
## 5 Concluding remarks
Using a quark potential model capable to describe consistently baryonic systems as three quark bound states we have derived an intrinsic transverse momentum distribution $`g(\stackrel{}{k}_T)`$ of partons inside the proton. This, clearly non-Gaussian, distribution is characterized by the presence of a smooth local maximum at relatively high transverse momenta. Our approach is based on the idea that transverse momentum effects may be described nonrelativisticaly as they are not influenced by the Lorenz boost along the beam axis. Describing the $`k_T`$-smearing effects in the parton model for $`pp`$ collisions through this non-Gaussian distribution we calculated the differential cross section for $`\pi ^0`$-production at midrapidity for different beam energies. Assuming that the corresponding nonperturbative parameter $`<k_T>`$, related to $`g(\stackrel{}{k}_T)`$, depends on $`p_T`$ of the finally produced hadron as well as the incident energy, we obtain a perfect description of the experimental data measured in $`pp`$ collisions at three different experiments. The corresponding values of $`<k_T>`$ as a function of $`p_T`$ could originate, according to Heisenberg uncertainty relation, from the internal partonic structure of the proton. Our approach is expected to be valid also in single photon production as well as $`pA`$ and $`AA`$ processes . Therefore the performed analysis shows that many-body effects through a confining potential, reflected at the level of one-particle distributions, may influence strongly the $`k_T`$-smearing phenomena observed in hadronic collisions and therefore should be taken into account for a correct description of the experimental data.
Acknowledgment We thank N.G. Antoniou for helpful discussions. This work is financed by EPEAEK in the framework of PYTHAGORAS grants supporting University research groups under contract 70/3/7420.
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# Extinction properties of the X-ray bright/optically faint afterglow of GRB 020405
## 1 Introduction
Long-duration Gamma Ray Bursts (GRBs) are associated with the core collapse of massive stars exploding as type Ic supernovae (GRB 980425, Galama et al. 1998; GRB 011211, Della Valle et al. 2003; GRB 030329, Stanek et al. 2003; GRB 031203, Malesani et al. 2004). Due to their short life-time, massive stars are likely to die within their birthplace. Consequently, long-duration GRBs occur in the same star formation region where their progenitor star was born and where it rapidly evolved. If high redshift galaxy star formation regions are akin to Galactic giant molecular clouds, a dense and dusty environment is expected in the vicinity of GRBs.
Indeed, the measured equivalent hydrogen column densities $`N_H`$ obtained from X-ray afterglow spectral analysis are on average consistent with those observed along the line of sight of Galactic molecular clouds (e.g. Galama & Wijers 2001, De Pasquale et al. 2003, Stratta et al. 2004). In some afterglows (e.g. GRB 971214, Ramaprakash et al. 1998; GRB 980703 Bloom et al. 1998; GRB 980329, Yost et al. 2002) a high visual extinction $`A_V`$ in the GRB host galaxy has been measured. In addition, a large dust content was inferred from high resolution spectroscopy by refractory metal abundance analysis along the line of sight of four optical afterglows (Savaglio et al. 2003, 2004). However, on average, the estimated afterglow reddening is low and the rest frame visual extinction is a factor 10-100 lower than expected from the extrapolation of the X-ray $`N_H`$, assuming the Galactic dust-to-gas ratio (ล imon et al. 2001; Galama & Wijers 2000). In a few high redshift GRBs for which high resolution afterglow spectroscopy was performed, high $`N_H`$ from $`Ly_\alpha `$ absorption was measured. On the other hand a low $`A_V`$ is inferred from continuum absorption modeled with the Small Magellanic Cloud (SMC) extinction law. These results can be explained by a low metallicity and/or a low dust-to-gas ratio in the burst environment (e.g. Hjorth et al. 2003, Vreeswijk et al. 2004).
Alternatively, the GRB environment may be characterized by a โnon-standardโ extinction law (neither Galactic nor SMC; Stratta et al. 2004, Savaglio et al. 2004). In particular, an extinction law weakly dependent on wavelength can provide high visual extinction without substantial reddening (e.g. Hjorth et al. 2003). Dust properties of high redshift environments - such as the GRB host galaxies - are still poorly known. We have some indications from the Magellanic Clouds (e.g. Pei et al. 1992) and from a sample of local starburst galaxies (e.g. Calzetti et al. 1997) and Active Galactic Nuclei (e.g. Maiolino, Marconi & Oliva 2001, Maiolino 2004) that the dust composition, grain size distribution and dust to gas ratio can be significantly different from the ones observed in our Galaxy. The GRB star-forming environment, in addition, may not be representative of the typical low-density host galaxy inter-stellar matter (ISM). Finally, the intense X-ray and UV flux from the burst itself can modify the intrinsic dust grain size distribution by differential dust grain destruction, since sublimation processes are more effective on small grains (e.g. Waxman & Draine 2000; Fruchter et al. 2001; Draine & Hao 2002; Perna & Lazzati 2002; Perna, Lazzati & Fiore 2003).
A diagnostic tool of the GRB environment extinction properties is the simultaneous optical-to-X-ray afterglow continuum spectral analysis. The intrinsic optical continuum can be extrapolated from the X-ray data and dust absorption can be measured from the deviations of the measurements from the expected fluxes. Good quality data are necessary to constrain the free parameters of the continuum model (the position of the cooling break possibly laying between the optical and X-ray bands). This method was applied, with a various degree of accuracy, in a handful of GRBs (Stratta et al. 2004).
In this work we present a study of the optically faint/X-ray bright afterglow of GRB 020405 (e.g. Dado et al. 2002, Masetti et al. 2003, Bersier et al. 2003, Covino et al. 2003, Price et al. 2003, Berger et al. 2003, Mirabal et al. 2003). This burst was discovered by the IPN on 2002 April 5.029 UT. The duration of the burst was $`40`$s, the 25-100 keV fluence was $`3\times 10^5`$ergs cm<sup>-2</sup> and the peak flux was $`10^6`$ ergs cm<sup>-2</sup> s<sup>-1</sup>. This GRB was also observed by the GRB Monitor on board BeppoSAX with a duration of $`60`$s in the 40-700 keV band and a 50-700 keV fluence of $`4\times 10^5`$ ergs cm<sup>-2</sup> s<sup>-1</sup> (e.g. Price et al. 2003). Optical observations started 18 hours after the burst event (Price et al. 2003) for a period of $`10`$ days. An unknown fading source was localized at R.A.13<sup>h</sup> 58<sup>m</sup> 03<sup>s</sup>.12 and Dec. -31 22 22<sup>โฒโฒ</sup>.2 with an uncertainty of 0<sup>โฒโฒ</sup>.3 (Masetti et al. 2002, Price et al. 2003). An exhaustive summary of all the afterglow observations performed in the optical band is presented in Masetti et al. (2003). The host galaxy has redshift $`z=0.691\pm 0.002`$ (Masetti et al. 2003). Chandra observations started on April 6.711 UT and lasted until 7.350 UT using LETGS in conjunction with the ACIS detector and revealed a new fading source that was identified as the X-ray afterglow counterpart. Mirabal et al. (2003) found a featureless spectrum, well described by a power law continuum with energy spectral index $`\alpha _X=0.72\pm 0.21`$ and a rest frame $`N_H`$ of $`(4.7\pm 3.7)\times 10^{21}`$cm<sup>-2</sup>. The X-ray light curve decayed as a power law, with a temporal decay index $`\delta _X=1.97\pm 1.10`$ (Mirabal et al. 2003). The optical (VRI) decay index is $`\delta _O=1.54\pm 0.06`$ from 1 to 10 days after the burst (Masetti et al. 2003). The NIR bands (J and H) show a less steep decay with index $`\delta _{IR}1.3`$ (Masetti et al. 2003). Radio observations show a rapidly fading โradio flareโ at 1.2 days from the burst event (Berger et al. 2003). Berger et al. (2003) found that the radio, optical and X-ray data are self-consistently modeled assuming a collimated ejecta expanding into a uniform medium. Best fit parameters yield a jet break time at about 1 day after the burst and indicate that the cooling frequency at that time is at energies lower than the optical band. The same conclusion on the location of the cooling frequency was obtained by Masetti et al. (2003), who showed how the optical spectral energy distribution has a break around the J band. For these reasons, for this burst there is convincing evidence for the absence of a spectral break between the X-ray and the optical band. The different decay index and spectral slope observed in the NIR band with respect to the optical may be due to the presence of the cooling frequency between the optical and the NIR bands (Masetti et al. 2003; Berger et al. 2003).
The paper is organized as follows: in ยง2 we summarize the data analysis procedure; in ยง3 we describe the adopted extinction curves; in ยง4 and ยง5 we present and discuss our results.
## 2 Data analysis procedure
We followed the standard procedure for the X-rays (0.1-10.0 keV) data reduction for Chandra ACIS data. We grouped energy channels in order to have at least 20 counts per bin in order to apply the $`\chi ^2`$ statistic in the fitting procedure. We first fit the X-rays spectrum with a power law model and two absorber components. The column density of the first absorber was fixed to the Galactic value toward the line of sight of the burst $`4.3\times 10^{20}`$cm<sup>-2</sup> (Dickey & Lockmann 1995). We subsequently allowed for the second absorber a free column density to constrain a possible contribution of absorption from the host galaxy ISM. We found that the data are well fitted by this model, with a best fit energy spectral index of $`\alpha _X=1.0\pm 0.2`$ and rest frame ($`z=0.691`$) hydrogen column density of $`N_{Hz}=(0.8\pm 0.2)\times 10^{22}`$cm<sup>-2</sup> in addition to the Galactic absorption ($`\chi ^2`$/d.o.f.=22.3/24, where d.o.f. stands for degrees of freedom). The best fit model (absorbed power law) 1.6-10.0 keV flux is $`8.7\times 10^{13}`$ ergs cm<sup>-2</sup> s<sup>-1</sup>. Errors are at 1 $`\sigma `$ level. We note that our spectral index is steeper than, but still consistent with, the one obtained by Mirabal et al. (2003).
The optical-NIR SED was obtained by taking the photometric points from the journal table published by Masetti et al. (2003). Observations in each photometric band (U,B,V,R,I,J,H and K) were performed at different times and with different telescopes. We extrapolated the magnitudes at the observational epoch $`t_0`$ of 7.01 UT (1.98 days after the burst). To this purpose, we firstly selected the observations performed as closer as possible to $`t_0`$, and then we extrapolated the magnitude at $`t_0`$ assuming a single powerlaw with decay index estimated at that time from Masetti et al. (2003). The decay index uncertainty was propagated to the magnitude errors with an average increase that depends on the temporal distance of the photometric measure from $`t_0`$ (see Tab.1 in Masetti et al. 2003). This results in an increase of the error on the magnitudes of 20% or more, depending on the temporal distance from the selected $`t_0`$. We choose $`t_0`$ as the average time of the X-ray observations given the large uncertainty in the X-ray decay slope (Mirabal et al. 2003).
We corrected the magnitudes for Galactic extinction (E(B-V)=0.054 mag) using the Galactic dust infrared map by Schlegel et al. (1998). Assuming a total-to-selective extinction $`R_V=A_V/E(BV)=3.1`$, we derived the extinction value in each photometric band from the Galactic extinction curve parameterization $`A(\lambda )/A_V`$ by Cardelli et al. (1989). We finally converted magnitudes to fluxes using the effective wavelengths and normalization fluxes given by Fukugita et al. (1995). A $`7\%`$ systematic error was added quadratically to the magnitude errors to account for mag-to-flux conversion uncertainties and for intercalibration errors from different telescopes (see also Masetti et al. 2001). The resulting magnitudes are summarized in Table 1.
A simultaneous fit of the NIR-to-X-rays spectrum was performed with the X-ray spectrum corrected for the total photoelectric absorption and the optical SED corrected for the Galactic extinction but not corrected for an extra-Galactic component. For the intrinsic SED, we considered both a single powerlaw and a broken powerlaw, the latter with break energy below the optical range. In order to estimate the host galaxy extinction we multiplied the continuum spectral model by an absorption component. We tested different ISM dust composition and dust to gas ratios by assuming several type of extinction curves (see ยง3).
## 3 Host galaxy extinction
Since we do not know a priori the host galaxy extinction curve, we tested several types of environments (Fig. 1) with different dust grain size distributions, dust compositions and dust to gas mass ratios. The ratio between the hydrogen column density $`N_H`$ and the visual extinction $`A_V`$ provides a measure of the dust-to-gas ratio for a given dust distribution and composition. For simulated time-dependent extinction laws, the dust-to-gas ratio at the relevant time is derived from the simulations.
We tested the following extinction curves:
* the Galactic extinction curve (hereafter G) from Cardelli et al. (1989), for which $`N_H/A_V=0.18\times 10^{22}`$cm<sup>-2</sup> (Predehl & Schmidt 1995).
* The Small Magellanic Cloud extinction curve (hereafter SMC) from Pei et al. (1992), for which $`N_H/A_V=1.6\times 10^{22}`$cm<sup>-2</sup> (Weingartner & Draine 2000).
* The attenuation curve derived for a sample of local starburst galaxies (hereafter C) from Calzetti et al. (1994). No $`N_H/A_V`$ has been measured for this case due to the complexity of the geometry of the dust and star distribution inside these galaxies (Calzetti et al. 2001).
* Two extinction curves obtained by Maiolino et al. (2001) for the environment of a sample of AGNs. The Q1 extinction curve is computed assuming a power law dust grain size distribution $`dn(a)a^qda`$ in the range $`a_{min}=0.005\mu <a<a_{max}=10\mu `$ and $`q=3.5`$. In this case $`N_H/A_V=0.7\times 10^{22}`$ cm<sup>-2</sup>. The Q2 extinction curve is derived assuming $`a_{min}=0.005\mu <a<a_{max}=1\mu `$ and $`q=2.5`$ and it yields $`N_H/A_V=0.3\times 10^{22}`$ cm<sup>-2</sup>.
* Extinction curves resulting from numerical simulations performed with the code of Perna & Lazzati (2002). This code computes the temporal evolution of the dust grain size distribution, taking into account grain erosion due to UV and X-ray illumination. The ionizing continuum is tailored to the case of GRB 020405 (see above). Only the 40 seconds of the burst emission were considered, since the overall fluence of the later afterglow was smaller than that of the prompt emission by at least an order of magnitude. Initial conditions were either a Galactic (final curve labelled G\_40) or SMC (final curve labelled SMC\_40) ISM. A uniform cloud with radius $`R=7\times 10^{22}`$ cm and density $`n=4\times 5`$ cm<sup>-3</sup> was considered. We checked that changing this set-up would not change our conclusions.
## 4 Results
We fit the NIR/optical and the X-ray spectral energy distribution, extrapolated to a common epoch (see ยง2), with a dust-absorbed power law model. The NIR-to-X-ray spectral index $`\alpha `$ was let free to vary. We found that the optical/NIR data corrected for the โstandardโ extinction curves, namely the G and the SMC, are incompatible with X-ray data (see Tab. 2). Unacceptable fits were obtained also with the Calzetti, the Q1 and the G\_40 curves (see ยง3).
Acceptable fits with comparable probability $`P(>\chi ^2)`$ of $`15\%`$, were instead obtained with the Q2 and the SMC\_40 curves. These curves have the weakest dependence on wavelengths among those considered (see Fig. 1). We note that the relatively high $`\chi ^2`$ levels are due to the apparent bend of the continuum in the NIR which cannot be reproduced by a power-law (Fig. 2).
We tested therefore a smoothed broken power law spectrum (with the sharpness parameter $`s`$ equal to 1, Granot & Sari 2002) with break energy between the optical and NIR bands. The optical-to-X-rays spectral index was let free to vary within the 90$`\%`$ confidence level interval 0.8-1.0, found from the X-ray analysis. Assuming $`\nu _c<\nu _X,\nu _O`$, from the X-ray spectral slope we derive an electron spectral index $`p=2.0\pm 0.4`$. We then compute the expected NIR spectral slope (corresponding to $`(p1)/2`$ for $`\nu _{NIR}<\nu _c`$, Sari et al. 1998) that we let free to vary within the 90$`\%`$ confidence level interval ($`0.30.7`$). Even allowing for this extra degree of freedom, we found that the G, SMC, C, Q1 and G\_40 absorption models are inconsistent with the data (see Tab. 3).
For the models Q2 and SMC\_40 we find an improvement in the fit with a $`2.5\sigma `$ statistical significance, according to the F-test (Fig. 3). They both fit the data successfully (Tab. 3). We then checked the consistency of the measured $`A_V`$ values with the $`N_{Hz}`$ measured from the X-ray analysis and the $`N_H/A_V`$ relationship expected (see ยง3). For the Q2 model, we found $`N_H=0.7\times 10^{22}`$cm<sup>-2</sup> assuming a power law model and $`N_H=0.8\times 10^{22}`$cm<sup>-2</sup> assuming a broken power law. These $`N_H`$ values are consistent with the $`N_{Hz}`$ \[($`0.8\pm 0.2)\times 10^{22}`$ cm<sup>-2</sup>\] measured from X-ray analysis. For the SMC\_40 model, even though the shape of the extinction curve fits the data well, a gray extinction with $`A_V2`$ is obtained only for $`N_H10^{24}`$ cm<sup>-2</sup>, largely exceeding the X-ray measured $`N_{Hz}`$. Formally a self consistent fit can be obtained assuming a dust-to-gas ratio 100 times larger than the SMC one. This would imply a $`10`$ times solar metallicity. We consider such physical conditions too extreme, even tough GRB (star formation) environments are expected to be especially dusty.
## 5 Discussion
We have studied the broad-band spectrum of the afterglow of GRB 020405 from NIR to X-ray bands. The optical-NIR and X-ray spectral indices are mutually consistent. However, this burst is peculiar since it has an optical-NIR flux a factor of $`10`$ times dimmer than the extrapolation of the X-ray spectrum. This can be attributed to an inverse Compton component in the X-ray band (Sari & Esin 2001). Alternatively, the optical brightness and spectrum could be affected by dust extinction. We investigated the latter possibility. We found a self-consistent model that allowed us to constrain the properties of the dust distribution and its formation history.
From a simultaneous NIR-to-X-ray spectral analysis, we found that the optical data can be fit with the external shock model only if corrected for an extinction curve weakly dependent on wavelength. Such an extinction curve can be the result of small dust grains coagulation and/or dust grain destruction (e.g. Maiolino et al. 2001, Perna et al. 2003). In Figure 1 we plot the extinction curves derived from these two different processes (see ยง3). The remarkable similarity of these extinction curves, computed assuming different processes, is evident.
The dust coagulation rate increases with density $`n^{1/2}`$ (Draine 1985). It is therefore favored in dense environments such as the cores of star formation regions where long GRBs are expected to explode. Incidentally, extinction curves weakly dependent on wavelengths have been inferred also for other extra-galactic objects, such as AGNs (Maiolino et al. 2001).
A dust distribution biased towards large grains can also be due to grain destruction mechanisms (Waxman & Draine 2000; Fruchter et al. 2001; Draine & Hao 2002; Perna & Lazzati 2002). We have attempted to reproduce the observations starting from known dust distributions in the local Universe (Galactic and SMC). The evolution of the initial distributions is followed as a result of the grain interaction with the burst ionizing flux (Perna et al. 2003). We find that the modified Galactic distribution cannot reproduce the observations, while the SMC one can reproduces them if an extremely large dust-to-gas ratio and metallicity are assumed. Since these conditions are hardly realized in nature, a scenario in which the large grain distribution is intrinsic to the host galaxy is favored.
The difficulty in reproducing the results as a consequence of dust destruction can be understood since dust destruction reduces the opacity besides modifying the reddening law. Since this burst requires a large extinction on a day timescale, it is likely that the opacity is provided by dust at a large distance from the burster, which is not affected by the burst photons. Prompt optical-NIR observations are better suited for the detection of dust destruction, since the evolution is caught โon the actโ. This would allow us to detect dust before it is destroyed and, based on the properties of what is left, infer the geometrical and physical conditions of the host galaxy ISM (Perna et al. 2003).
###### Acknowledgements.
The authors thank the referee F. J. Castander for his useful comments on the manuscript. G.S. is supported by the Research Training Network *Gamma-Ray Bursts: an Enigma and a Tool* funded by the EU.
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# Coherence Phenomena in the Phase-sensitive Optical Parametric Amplification inside a Cavity
## Abstract
We theoretically and experimentally demonstrate coherence phenomena in optical parametric amplification inside a cavity. The mode splitting in transmission spectra of phase-sensitive optical parametric amplifier is observed. Especially, we show a very narrow dip and peak, which are the shape of $`\delta `$ function, appear in the transmission profile. The origin of the coherence phenomenon in this system is the interference between the harmonic pump field and the subharmonic seed field in cooperation with dissipation of the cavity.
$`Introduction`$. โ Coherence and interference effects play very important roles in determining the optical properties of quantum systems. Electromagnetically induced transparency (EIT) one in quantum-mechanical atomic systems is a well understood and thoroughly studied subject. EIT has been utilized in a variety of applications, such as lasing without inversion two , slow and stored light two1 ; two2 , enhanced nonlinear optics three , and quantum computation and communication four . Relying on destructive quantum interference, EIT is a phenomenon where the absorption of a probe laser field resonant with an atomic transition is reduced or even eliminated by the application of a strongly driving laser to an adjacent transition. Since EIT results from destructive quantum interference, it has been recently recognized that similar coherence and interference effects also occur in classical systems, such as plasma five , coupled optical rsonatorssix , mechanical or electric oscillators seven . In particular, the phenomenology of the EIT and dynamic Stark effect is studied theoretically in a dissipative system composed by two coupled oscillators under linear and parametric amplification using quantum optics model in Ref. eight . The classical analog of EIT is not only helpful to understand deeply the physical meaning of EIT phenomenon, but also offers a number of itself important applications, such as slow and stored light by coupled optical resonatorsnine .
In this Letter, we extend the model in Ref. eight and present a new system - phase-sensitive optical parametric amplifier (OPA) to demonstrate coherence effects theoretically and experimentally. We observe mode splitting in transmission spectra of OPA. Especially, we show a very narrow dip and peak, which are the shape of $`\delta `$ function, appear in the transmission profile. This phenomenon results from the interference between the harmonic pump field and the subharmonic seed field in OPA. The destructive and constructive interference correspond to optical parametric deamplifier and amplifier respectively, which are in cooperation with dissipation of the cavity. The absorptive and dispersive response of an optical cavity for the probe field is changed by optical parametric interaction in the cavity. Phase-sensitive optical parametric amplifier presents a number of new characteristics of coherence effects.
$`Theoretical`$ $`model.`$ โ Consider the interaction of two optical fields of frequencies $`\omega `$ and $`2\omega `$, denoted by subharmonic and harmonic wave (the pump), which are coupled by a second-order, type-I nonlinear crystal in a optical cavity as shown in Fig.1. The cavity is assumed to be a standing wave cavity, and only resonant for the subharmonic field with dual-port of transmission $`T_{HR}`$ and $`T_c`$, internal losses $`A`$ and length $`L`$ (roundtrip time $`\tau =2L/c`$). We consider both the subharmonic seed beam $`a^{in}`$ and harmonic pump beam $`\beta ^{in}`$ are injected into the back port ($`T_{HR}`$ mirror) of the cavity, where the relative phase between the injected field is adjusted by a movable mirror outside the cavity. $`T_{HR}`$ mirror is a high reflectivity mirror at the subharmonic wavelength, yet has a high transmission coefficient at the harmonic wavelength and $`T_c`$ mirror has a high reflectivity coefficient for the harmonic wave. The harmonic wave makes a double pass through the nonlinear medium. The equation of motion for the mean value of the subharmonic intra-cavity field can then be derived by the semiclassical method ten as
$$\tau \frac{da}{dt}=i\tau \mathrm{\Delta }a\gamma a+g\beta ^{in}a^{}+\sqrt{2\gamma _{in}}a^{in}.$$
(1)
The decay rate for internal losses is $`\gamma _l=A/2`$ and the damping associated with coupling mirror and back mirror is $`\gamma _c=T_c/2`$ and $`\gamma _{in}=T_{HR}/2`$, respectively. The total damping is denoted by $`\gamma =\gamma _{in}+\gamma _c+\gamma _l`$. $`\mathrm{\Delta }`$ is the detuning between the cavity-resonance frequency $`\omega _c`$ and the subharmonic field frequency $`\omega `$. The strength of the interaction is characterized by the nonlinear coupling parameter $`g`$. Eq.1 is complemented with the boundary conditions $`a^{out}=\sqrt{2\gamma _c}a`$ and $`a^{ref}=a^{in}+\sqrt{2\gamma _{in}}a`$ to create propagating beams, where $`a^{out}`$ is the transmitted field from the coupling mirror $`T_c`$ and $`a^{ref}`$ is the reflected field from the back mirror $`T_{HR}`$. The phase-sensitive optical parametric amplifier always operates below the threshold of optical parametric oscillation (OPO) $`\beta _{th}^{in}=\gamma /g`$. Eq.1 ignores the third-order termeleven describing the conversion losses due to harmonic generation. For simplicity, we assume that the phase of the pump field is zero in any case, i.e, $`\beta ^{in}`$ is real and positive value. The intra-cavity field $`a`$ and the injected field $`a^{in}`$ are expressed as $`a=`$ $`\alpha \mathrm{exp}\left(i\varphi \right)`$ and $`a^{in}=`$ $`A_{in}\mathrm{exp}\left(i\phi \right)`$ respectively. Here, $`\alpha `$ and $`A_{in}`$ are real values, $`\varphi `$ and $`\phi `$ are the relative phase between the intra-cavity field and the pump field as well as between the seed field and the pump field, respectively. If the harmonic pump is turned off, the throughput for the non-impedance matched subharmonic seed beam is given by $`a_{nopump}^{out}=2\sqrt{\gamma _c\gamma _{in}}A_{in}/(\gamma +i\tau \mathrm{\Delta })`$. The subharmonic seed beam is subjected to either amplification or de-amplification, depending on the chosen relative phase between the subharmonic field and the pump field.
$`Case1:`$ Consider the transmitted intensity of the subharmonic seed beam as a function of the detuning $`\mathrm{\Delta }`$ between the subharmonic field frequency and the cavity-resonance frequency, and keep the pump field of frequency $`\omega _p=2\omega `$ constant. Setting the derivative to zero ($`d\alpha /dt=0`$) and separating the real and image part of Eq.1, the steady state solutions of the amplitude and relative phase of the intra-cavity field are given by
$`\gamma \alpha +g\beta ^{in}\alpha \mathrm{cos}2\varphi +\sqrt{2\gamma _{in}}A_{in}\mathrm{cos}\left(\varphi \phi \right)`$ $`=`$ $`0,`$ (2)
$`\tau \mathrm{\Delta }\alpha +g\beta ^{in}\alpha \mathrm{sin}2\varphi +\sqrt{2\gamma _{in}}A_{in}\mathrm{sin}\left(\varphi \phi \right)`$ $`=`$ $`0.`$
When the amplitude and relative phase of the subharmonic seed beam are given, the transmitted intensity of the subharmonic beam is obtained from Eq.2 and the boundary condition. Fig.2(a) shows a Lorentzian profile of the subharmonic transmission when the pump field is absent. This corresponds to the typical transmitted spectrum of the optic al empty cavity. When the injected subharmonic field is out of phase ($`\phi =\pi /2)`$ with the pump field, the subharmonic transmission profile is shown in Fig.2(b,c,d) for different pump powers, in which there is a symmetric mode splitting. The transmitted power of the subharmonic beam is normalized to the power in the absence of the pump and zero detuning. The transmission spectra show that the dip becomes deeper and two peaks higher as the pump intensity increases. The origin of mode splitting in transmission spectra of OPA is destructive interference in cooperation with dissipation of the cavity. If the subharmonic field resonates in the cavity perfectly, i.e. $`\mathrm{\Delta }=0`$, the subharmonic intra-cavity field and the pump filed are exactly out of phase and will interfere destructively to produce the deamplification for the subharmonic field in the nonlinear crystal. Thus a dip appears at the zero detuning of the transmission profile. If the subharmonic field is not quite resonant in the cavity perfectly, that is, the subharmonic fieldโs frequency is not exactly an integer multiple of the free spectral range (but close enough to build up a standing wave), the phase difference between the subharmonic intra-cavity field and the pump field will not be exactly out of phase and will increase as the detuning increases. The subharmonic intra-cavity field will change from deamplification to amplification as the phase difference increases. Thus we see that the transmission profile has two symmetric peaks at two detuning frequencies. When the phase of the injected subharmonic field is deviated from out of phase with the pump field, i.e. $`\phi =\pi /2\pm \theta `$, an asymmetric mode splitting in the subharmonic transmission profile is illustrated in Fig.2(e,f), in which the dip is deviated from the zero detuning of the transmission profile and two peaks have different amplitude.
$`Case2:`$ Consider the subharmonic transmission profiles when the frequency of the pump field is fixed at $`\omega _p=2(\omega _c+\mathrm{\Omega })`$. When scanning the frequency of the the subharmonic seed beam, an idler field in the OPA cavity will be generated with the frequency $`\omega _i=\omega _p\omega `$ due to energy conservation. The equation of motion of OPA become frequency-nondegenerate and is given by
$`\tau {\displaystyle \frac{da}{dt}}`$ $`=`$ $`i\tau \mathrm{\Delta }a\gamma a+g\beta ^{in}a_i^{}+\sqrt{2\gamma _{in}}a^{in},`$ (3)
$`\tau {\displaystyle \frac{da_i}{dt}}`$ $`=`$ $`i\tau \mathrm{\Delta }_ia_i\gamma a_i+g\beta ^{in}a^{}`$
where $`a_i`$ is the idler field in the OPA cavity. $`\mathrm{\Delta }_i`$ is the detuning between the cavity-resonance frequency $`\omega _c`$ and the idler field frequency $`\omega _i`$. Thus the subharmonic transmission profile in this case is obtained from Eq.3 for $`\omega \omega _i`$ and Eq.1 for $`\omega =\omega _i`$. When $`\mathrm{\Omega }=0`$, so $`\mathrm{\Delta }=\mathrm{\Delta }_i`$, the stationary solution of the subharmonic and idle field is given by solving the mean-field equations of Eq.3 and using the input-output formalisms. We obtain
$`A^{out}`$ $`=`$ $`{\displaystyle \frac{2\sqrt{\gamma _c\gamma _{in}}}{i\tau \mathrm{\Delta }+\gamma \frac{(g\beta ^{in})^2}{i\tau \mathrm{\Delta }+\gamma }}}A^{in},`$ (4)
$`A_i^{out}`$ $`=`$ $`{\displaystyle \frac{2\sqrt{\gamma _c\gamma _{in}}g\beta ^{in}}{\left(i\tau \mathrm{\Delta }+\gamma \right)^2(g\beta ^{in})^2}}A^{in}.`$
We will record the total output power including the subharmonic and idle field. The transmitted power of the subharmonic beam is given by
$$P_{out}^{nor}=\{\begin{array}{c}\left|\frac{\gamma }{i\tau \mathrm{\Delta }+\gamma \frac{(g\beta ^{in})^2}{i\tau \mathrm{\Delta }+\gamma }}\right|^2+\left|\frac{\gamma g\beta ^{in}}{\left(i\tau \mathrm{\Delta }+\gamma \right)^2(g\beta ^{in})^2}\right|^2\\ \mathrm{if}\omega \omega _i\\ \frac{\gamma ^2}{(\gamma \pm g\beta ^{in})^2}\mathrm{if}\omega =\omega _i.\end{array}$$
(5)
Here, $`\pm `$ corresponds to the deamplifier and amplifier in frequency-degenerate OPA. Fig.3(a) and (b) show that the very narrow dip and peak, which is the shape of $`\delta `$ function, appear in the transmission profile. This novel coherence phenomena results in that the destructive and constructive interference are established only in the point of $`\omega =\omega _i`$, and completely destroyed in the other frequencies.
$`Experiment`$. โ The experimental setup is shown schematically in Fig.4. A diode-pumped intracavity frequency-doubled continuous-wave(cw) ring Nd:YVO$`_\text{4}`$/KTP single-frequency green laser severs as the light sources of the pump wave (the second-harmonic wave at $`532`$ $`nm`$) and the seed wave (the fundamental wave at $`1064`$ $`nm`$) for OPA. The green beam doubly passes the acousto-optic modulator (AOM) to shift the frequency 440 MHz. The infrared beam doubly passes AOM to shift the frequency around 220 MHz. We actively control the relative phase between the subharmonic and the pump field by adjusting the phase of the subharmonic beam with a mirror mounted upon a piezoelectric transducer (PZT). Both beams are combined together by a dichroic mirror and injected into the OPA cavity. OPA consists of periodically poled KTiOPO<sub>4</sub> (PPKTP) crystal (12 $`mm`$ long) and two external mirrors separated by $`63`$ $`mm`$. Both end faces of crystal are polished and coated with an antireflector for both wavelengths. The crystal is mounted in a copper block, whose temperature was actively controlled at millidegrees kelvin level around the temperature for optical parametric process (31.3C). The input coupler M1 is a $`30`$ $`mm`$ radius-of-curvature mirror with a power reflectivity $`99.8\%`$ for $`1064`$ $`nm`$ in the concave and a total transmissivity $`70\%`$ for $`532`$ $`nm`$, which is mounted upon a PZT to adjust the cavity length. The output wave is extracted from M2, which is a $`30mm`$ radius-of-curvature mirror with a total transmissivity $`3.3\%`$ for $`1064`$ $`nm`$ and a reflectivity $`99\%`$ for $`532`$ $`nm`$ in the concave. Due to the large transmission of input coupler at $`532`$ $`nm`$, the pump field can be thought as only passes the cavity twice without resonation. The measured cavity finesse was $`148`$ with the PPKTP crystal, which indicates the total cavity loss of $`4.24`$%. Due to the high nonlinear coefficient of PPKTP, the measured threshold power is only $`35`$ $`mW`$.
First, we fix the frequency of the subharmonic and the pump field with $`\omega _p=2\omega `$ and scan cavity length, which corresponds to the condition of case 1. Figure 5 shows the experimental results: (a) without the pump field, (b) $`\phi =\pi /2`$ and $`\beta ^{in}/\beta _{th}^{in}=0.33`$, (c) $`\phi =\pi /2`$ and $`\beta ^{in}/\beta _{th}^{in}=0.71`$, (d) $`\phi =\pi /2`$ and $`\beta ^{in}/\beta _{th}^{in}=0.9`$, (e) $`\phi =\pi /20.07`$ and $`\beta ^{in}/\beta _{th}^{in}=0.9`$, (f) $`\phi =\pi /2+0.07`$ and $`\beta ^{in}/\beta _{th}^{in}=0.9`$. It can be seen that the experimental curves are in good agreement with the theoretical results shown in Fig.2, which are obtained with the experimental parameters.
Then, we fix the cavity length and frequency of the pump field and scan the frequency of the subharmonic field by the AOM, which corresponds to the condition of case 2. The output including the subharmonic and idle field is detected by a photodiode. There is a beat-note signal in the photocurrent with frequency proportional to the detuning. The very narrow dip and peak appeared in a broad Lorentzian profile are observed experimentally as shown in Fig.6. The insets in Fig.6 show the enlarged narrow dip and peak by reducing the scanned range of frequency, which present the square shape. Because the measurement of transmission profile is dynamic processes, the shape of $`\delta `$ function for the narrow dip and peak in the theoretical model becomes square shape in experiment. The width of the square shape is $``$2KHz which is estimated from the voltage on VCO (Voltage-Controlled Oscillator) of AOM.
$`Conclusion.`$ โ We reported the theoretical and experimental results of coherence phenomena in the phase-sensitive optical parametric amplification inside a cavity. The splitting in transmission spectra of OPA was observed. Mode splitting, as well known, occurs not only in coupled quantum system, but also in coupled optical resonators and in coupled mechanical and electronic oscillators. To the best of our knowledge, we first observed mode splitting experimentally in the optical parametric process. This system will be important for practically optical and photonic applications such as optical filters, delay lines, and closely relate to the coherent phenomenon of EIT predicted for quantum systems. OPA also has a important application as squeezed light source. Our results may help us to investigate quantum noise spectrum.
Corresponding authorโs email address: jzhang74@yahoo.com, jzhang74@sxu.edu.cn
###### Acknowledgements.
J. Zhang thanks Prof. Kunchi Peng and Changde Xie for the helpful discussions.This research was supported in part by National Natural Science Foundation of China (Approval No.60178012), Program for New Century Excellent Talents in University, Natural Science Foundation of Shanxi Province, and the Research Fund for the Returned Abroad Scholars of Shanxi Province. REFERENCES
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# A Tractable Example of Perturbation Theory with a Field Cutoff: the Anharmonic Oscillator
## 1 Introduction
Perturbative methods and Feynman diagrams have played an important role in the development of quantum field theory and its applications. However, perturbative series usually have a zero radius of convergence. For scalar models with $`\lambda \varphi ^4`$ interactions, the coefficients of perturbative series grow factorially. For any fixed, strictly positive, value of $`\lambda `$, there exists an order beyond which adding higher order terms diminishes the accuracy. This feature will restrict our ability to perform high precision tests of the standard model (for instance, $`g2`$ of leptons and the hadronic width of the $`Z^0`$) during the next decades.
The large order behavior of the series is dominated by large field configurations which have little effect on low energy observables. Introducing a large field cutoff in the path integral formulation of scalar field theory, dramatically improves the large order behavior of the perturbative series. In two non-trivial examples , this procedure yields series that have finite radii of convergence and tend to values that are exponentially close to the exact ones. This also allows us to define the theory for negative or complex values of $`\lambda `$, a subject that has raised a lot of interest recently . An important feature of this approach is that for a perturbative expansion at a given order in $`\lambda `$, it is possible in some cases to determine an optimal field cutoff using the strong coupling expansion , bridging the gap between the two expansions. In other words, the modified perturbative methods allow us to take into account non-perturbative effects.
Despite these promising features, calculating the modified coefficients remains a challenging technical problem. While developing a new perturbative method (see e.g. Ref. for the $`\delta `$ expansion), it is customary to demonstrate the advantages of a method with simple integrals and the non-trivial, but well-studied, case of the anharmonic oscillator. A simple integral has been discussed in Ref. . In this article, we show not only that this program can be completed in the case of the anharmonic oscillator, but also that the results show remarkable properties:
* For the two lowest orders, the approximate formulas obtained in the large field cutoff limit extend unexpectedly far in the low field cutoff region.
* For the higher orders, the transition between the small field cutoff regime and the large field cutoff regime can be approximately described in terms of a single function.
The results are presented in the following way. In Sec. 2, we define the model considered and the numerical calculation of the modified coefficients. In Sec. 3, we discuss the radius of convergence of the modified series from a numerical point of view. In Secs. 4 and 5, we present the approximate methods at small and large field cutoff. Rigorous bounds on the radius of convergence of series encountered in the previous sections are given in Sec. 6. The question of the interpolation between the small and large field cutoff regimes is discussed in Sec. 7. The importance of a complete understanding of this question, as well as the extension to field theory are discussed in the conclusions.
## 2 The problem and its numerical solution
In this section, we introduce the anharmonic oscillator with a โfield cutoffโ and we explain how to calculate numerically the perturbative series in the anharmonic coupling $`\lambda `$. This method was first used in and tested with the known results up to order 20. It is a perturbative version of the numerical method proposed in Ref. . For convenience, we use quantum mechanical notations instead of field theoretical ones $`\varphi x,m\omega `$ and the field cutoff will be denoted $`x_{nax}`$. We also use units such that $`\mathrm{}=1`$ and the โmechanical massโ $`m`$ is 1. However, the harmonic angular frequency $`\omega `$, will sometimes be used as an expansion parameter and will be kept arbitrary in the equations. In these units, $`x_{max}\sqrt{\omega }`$ and $`\lambda /\omega ^3`$ are dimensionless. The hamiltonian reads
$$H=\frac{p^2}{2}+V(x)$$
(1)
with
$$V(x)=\{\begin{array}{ccc}\frac{1}{2}\omega ^2x^2+\lambda x^4& \mathrm{if}& |x|<x_{\mathrm{max}}\\ \mathrm{}& \mathrm{if}& |x|x_{\mathrm{max}}\end{array}$$
(2)
Our main goal is the calculation of the modified coefficients $`E_0^{(k)}(\sqrt{\omega }x_{max})`$ of the perturbative series for the ground state energy:
$$E_0(x_{max},\omega ,\lambda )=\omega \underset{k=0}{\overset{\mathrm{}}{}}E_0^{(k)}(\sqrt{\omega }x_{max})\times (\lambda /\omega ^3)^k,$$
(3)
For this purpose, we will solve perturbatively the time independent Schrรถdinger equation with the boundary condition $`\mathrm{\Psi }(x_{max})=0`$. We proceed as in Ref. . Setting
$$\mathrm{\Psi }(x)\mathrm{e}^{_{x_0}^x๐y(L(y)/K(y))},$$
(4)
the Schrรถdinger equation reads
$`L^{}`$ $`+`$ $`2(VE)K+GL=0`$ (5)
$`K^{}`$ $`+`$ $`L+GK=0`$ (6)
where $`G(x)`$ is an arbitrary function. The second of these equations implies that
$$\mathrm{\Psi }(x)K(x)\mathrm{e}^{^x๐yG(y)}$$
(7)
It possible to show that if $`G`$ and $`V`$ are polynomials:
* Eqs. (5-6) can be solved by power series which define entire functions; recursion formulas and initial conditions (different for even and odd solutions) are given in Ref. .
* the zeroes of $`\mathrm{\Psi }`$ are the sames as the zeroes of $`K`$.
This implies that the energy levels can be obtained by solving
$$K(x_{max},\omega ,\lambda ,E)=0$$
(8)
for the variable $`E`$, and that polynomial approximations can be used for this purpose. We now use the perturbative expansion
$$E=\omega \underset{k=0}{\overset{\mathrm{}}{}}E^{(k)}(\lambda /\omega ^3)^k.$$
(9)
We assume that $`\omega >0`$. The case $`\omega =0`$ is simpler and will be discussed in Sec. 5. By construction, the coefficients $`E^{(k)}`$ are dimensionless and can only depend on $`\sqrt{\omega }x_{max}`$ when their value is fixed using Eq. (8). In the following, this dependence will sometimes be kept implicit in order to reduce the size of some equations. We also recall that if we had not chosen the units $`\mathrm{}=m=1`$, the dimensionless quantities used above would read $`(\sqrt{\omega m/\mathrm{}}x_{max})`$ and $`\mathrm{}\lambda /m^2\omega ^3`$.
We need to set the expansion
$`K(x_{max},\omega ,\lambda ,E)`$ $`=`$ $`K^{(0)}(\sqrt{\omega }x_{max},E^{(0)})+`$ (10)
$`(\lambda /\omega ^3)K^{(1)}(\sqrt{\omega }x_{max},E^{(0)},E^{(1)})+\mathrm{}`$ (11)
equal to zero order by order in $`\lambda `$. We can approximate the $`K^{(k)}`$ by polynomials. We then start by solving $`K^{(0)}(\sqrt{\omega }x_{max},E^{(0)})=0`$ for $`E^{(0)}`$ using Newtonโs method. This correponds to the problem of an harmonic oscillator potential that becomes infinite at $`x=\pm x_{max}`$. The various zeroes of the polynomial correspond to the (even or odd depending on how we construct $`K`$ ) spectrum of this model. The ground state energy is obtained by taking the lowest even solution. By increasing the degree of the polynomial approximation, we can stabilize the numerical value with great accuracy. The independence on $`G`$ of the exact equations, before we use polynomial approximations, can be used to test the numerical accuracy of the polynomial approximations used for $`K`$. In practice a good choice is made by having $`K`$ of order 1 for $`E`$ close, but not fine tuned, to its actual value. Since the potential is even, it is natural to have $`K`$ even and $`G`$ and $`L`$ odd, however we noticed that introducing a small parity breaking in $`G`$ usually improves the numerical stability. We then solve $`K^{(1)}(x_{max},E^{(0)},E^{(1)})=0`$ for $`E^{(1)}`$. Since $`E^{(1)}`$ appears only linearly in the order $`\lambda `$ expansion of $`K`$, it is a linear equation for this quantity and it can be solved trivially. The same reasoning shows that the higher order equations are linear in $`E^{(k)}`$.
By using this method, we have calculated the first ten coefficients for $`\omega =1`$ and values of $`x_{max}`$ between $`0.5`$ and 7. In order to allow a comparison among the different orders, we define the ratios
$$R_k(\sqrt{\omega }x_{max})E_0^{(k)}(\sqrt{\omega }x_{max})/E_0^{(k)}(\mathrm{}),$$
(12)
which all tend to 1 in the $`x_{max}\mathrm{}`$ limit and to 0 in the $`x_{max}0`$ limit. The numerical values of these ratios are shown in Fig. 1. A striking feature is that the curves for the various orders have approximately the same shape and that we could approximately superpose them by appropriate horizontal translations. Before studying these curves in the large and small $`x_{max}`$ limits, we will discuss numerical estimates for the radius of convergence of the modified series, that can be extracted from this data.
## 3 Numerical evidence for a finite radius of convergence
If we calculate the integral of a function defined by a series which is uniformly convergent over the range of integration, it is legitimate to interchange the sum and the integral. On general grounds, one would thus expect that if we calculate the partition function of a properly regularized theory (say on a finite lattice) with a field cut, the perturbative expansion becomes convergent for arbitrary complex values of the coupling. The ground state energy of the anharmonic can be obtained from the logarithm of the partition function and consequently, we would expect that its perturbative expansion will have a finite radius of convergence. This radius will depend on the position of the complex zeroes of the partition function.
It is sometimes possible to estimate the radius of convergence by considering the empirical asymptotic behavior of the perturbative series. However, with a limited series, transient behavior is often encountered. It is clear from Fig. 1, that if $`x_{max}`$ is large enough, the beginning of the series looks like its $`x_{max}\mathrm{}`$ limit which grows at a factorial rate. Consequently, in order to attempt to probe the asymptotic behavior, we will only consider $`x_{max}`$ such that the available coefficients are significantly smaller than their value when $`x_{max}`$ is infinite. For this reason, we have limited the range of investigation to $`x_{max}<1.3`$. If we have a finite radius of convergence $`\lambda _c`$, we expect that for $`k`$ sufficiently large,
$$|E_0^{(k)}(\sqrt{\omega }x_{max})|(P(\sqrt{\omega }x_{max}))^k$$
(13)
When this is the case, $`\lambda _c(x_{max},\omega )=\omega ^3/P(\sqrt{\omega }x_{max})`$. In Fig. 3, we have plotted ln$`|E_0^{(k)}(\sqrt{\omega }x_{max})|`$ versus $`k`$. We see a clear linear behavior and the slope can be interpretated as ln$`P(\sqrt{\omega }x_{max})`$. We can then study $`\lambda _c(x_{max},\omega )`$ as a function of $`x_{max}`$. This is done in Fig. 3 where a log-log plot shows a clear linear behavior. The slope is near 5.9 which is close to the value 6 that, we will argue in Sec. 6, should hold in the limit of small $`x_{max}`$ (see Eq. (40)). We conclude that in this limit, our numerical data suggests
$$\lambda _c(x_{max},\omega =1)65\times x_{max}^6.$$
(14)
## 4 The large $`x_{\mathrm{max}}`$ limit
In this section, we first discuss the harmonic energy spectrum with vanishing boundary conditions at $`\pm x_{max}`$ and then treat the anharmonic interactions with the usual perturbative methods. We work in the limit where $`x_{max}`$ is large. This means that if we consider the $`n`$-th energy level, $`x_{max}`$ should be much larger than the largest zero of $`H_n(x\sqrt{\omega })`$. According to Eq. (6.32.5) of Ref. , this imposes the restriction $`x_{max}>>\sqrt{2n/\omega }`$ for large $`n`$. When all the zeroes of $`H_n`$ are well within $`[x_{max},x_{max}]`$, the new boundary conditions require changes that are exponentially small. It is clear that for any fixed $`x_{max}`$, this condition will be violated for $`n`$ sufficiently large.
### 4.1 The harmonic case
We first need to calculate the energy eigenvalues at $`\lambda =0`$, $`E_n^{(0)}(\sqrt{\omega }x_{max}^2`$), and their corresponding wave functions $`\mathrm{\Psi }_n^{(0)}(x)`$. The wave function always depends on $`x_{max}`$, but this will be kept implicit. We impose the boundary conditions $`\mathrm{\Psi }_n^{(0)^{}}(0)=0`$ for $`n`$ even, $`\mathrm{\Psi }_n^{(0)}(0)=0`$ for $`n`$ odd, and $`\mathrm{\Psi }^{(0)}(x_{max})=0`$ in both cases. We will not pay attention to the overall normalization until the end of the calculation. When $`x_{\mathrm{max}}\mathrm{}`$, we have $`E_n^{(0)}(\sqrt{\omega }x_{max})(n+1/2)`$ and we will use
$$ฯต_n(\sqrt{\omega }x_{max})E_n^{(0)}(\sqrt{\omega }x_{max})n1/2$$
(15)
as our expansion parameter. Again, $`ฯต_n`$ is a dimensionless quantity that can only depend on $`\sqrt{\omega }x_{max}`$ and we will often keep this dependence implicit. We write
$$\mathrm{\Psi }_n^{(0)}(x)K_n^{(0)}(x)\mathrm{e}^{\omega x^2/2}.$$
(16)
This corresponds to a choice $`G=\omega x`$ in Eqs. (5-6) and consequently
$$(1/2)K_n^{(0)^{\prime \prime }}+\omega xK_n^{(0)^{}}n\omega K_n^{(0)}=ฯต_n\omega K_n^{(0)}.$$
(17)
We then expand
$$K_n^{(0)}=K_n^{(0)(0)}+ฯต_nK_n^{(0)(1)}+\mathrm{}$$
(18)
When $`x_{max}\mathrm{}`$, $`ฯต_n0`$, and we want to recover the usual harmonic oscillator solution. Consequently, we set $`K_n^{(0)(0)}(x)=H_n(\sqrt{\omega }x)`$, the $`n`$-th Hermite polynomial which is obviously a solution of Eq. (17) at order 0 in $`ฯต_n`$. At order $`ฯต_n`$, we have
$$(1/2)K_n^{(0)(1)^{\prime \prime }}+\omega xK_n^{(0)(1)^{}}n\omega K_n^{(0)(1)}=\omega H_n(\sqrt{\omega }x),$$
(19)
with the boundary conditions
$$K_n^{(0)(1)}(0)=K_n^{(0)(1)^{}}(0)=0.$$
(20)
Remarkably, this inhomogeneous equation can be integrated exactly in two steps. First, we set
$$K_n^{(0)(1)}(x)=H_n(\sqrt{\omega }x)G_n(x).$$
(21)
This removes the explicitly $`n`$-dependent term and the equation depend only on $`G^{}`$:
$$(1/2)H_n(\sqrt{\omega }x)G_n^{^{\prime \prime }}+\left[\omega xH_n(\sqrt{\omega }x)H_n^{}(\sqrt{\omega }x)\sqrt{\omega }\right]G_n^{^{}}=\omega H_n(\sqrt{\omega }x),$$
(22)
We then write $`G^{}(x)=N(x)\mathrm{e}^{\omega x^2}`$ which removes the term linear in $`x`$ and allows us to write the l. h. s. as a total derivative. The solution is then
$$G_n(x)=2\omega _0^x๐y(H_n(\sqrt{\omega }y))^2\mathrm{e}^{\omega y^2}_0^y๐z\mathrm{e}^{\omega z^2}(H_n(\sqrt{\omega }z))^2,$$
(23)
One can check that the lower bounds of integration imply the boundary conditions of Eq. (20). This is obvious when $`n`$ is even. When $`n`$ is odd, the $`z`$ integral is of order $`y^3`$ when $`y0`$ which overcompensates the $`y^2`$ of the other factors. Note that $`G_n(x)`$ is independent of $`x_{max}`$. The integrand of the $`y`$-integral has a double pole at the zeroes of $`H_n`$, however it has no single pole. This can be seen by plugging a Laurent expansion of $`G^{}`$ with poles of order 1 and 2 about a zero of $`H_n`$, and using the relation between the first and the second derivative of $`H_n`$ at this zero provided by the defining equation for Hermite polynomials. This forces the coefficient of the simple pole to be zero. Consequently, when we express $`G`$ as the integral of $`G^{}`$, we can go around the double pole either above or below the real line and obtain the same result. In other words, we can regularize the $`y`$-integral by replacing $`(H_n(\sqrt{\omega }y))^2`$ by $`(H_n(\sqrt{\omega }y)\pm iฯต)^2`$. From these considerations, we see that $`G(x)`$ develops a simple pole when $`x`$ approaches a zero of $`H_n`$, which compensate the simple zero of $`H_n`$. Due to the absence of single pole in $`G^{}`$, there is no logarithmic singularity in $`G`$.
$`ฯต_n`$ can be calculated by imposing the condition that the wave function vanishes at $`x_{max}`$. This translates into the simple equation:
$$ฯต_n(\sqrt{\omega }x_{max})=1/G_n(x_{max}).$$
(24)
As $`ฯต_n`$ increases, the non-trivial zeroes of the wavefunction move toward the origin. The shift of these zeroes is approximately linear in $`ฯต_n`$ for large $`x_{max}`$. For instance, for $`n`$ =2, the shift of the non-trivial zero $`x_{(0)}`$ obeys the approximate linear relation $`\mathrm{\Delta }x_{(0)}0.145ฯต_2`$.
We now discuss in more detail the case of the ground state ($`n=0`$). It should be noted that in this case the double integral can be expressed in terms of generalized hypergeometric series
$`G_0(x)`$ $`=`$ $`\omega x_2^2F_2(1,1;2,3/2;\omega x^2)`$
$`=`$ $`{\displaystyle \frac{\sqrt{\pi }}{2}}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(\omega x^2)^k}{k\mathrm{\Gamma }(k+1/2)}}.`$
This series defines a function that converges over the entire complex plane. The asymptotic form of $`G_0(x_{max})`$, for $`x_{max}`$ large, can be worked out by noticing that the integral over $`y`$ in Eq. (23) comes mostly from the region $`yx_{max}`$ and consequently, we can extend the integral over $`z`$ to infinity with exponentially small errors. The integral over $`y`$ can then be approximately performed by expanding the argument of the exponential about $`x_{max}`$. Using Eq. (24), we conclude that asymptotically
$$ฯต_0(\sqrt{\omega }x_{max})2\sqrt{\frac{\omega }{\pi }}x_{\mathrm{max}}e^{\omega x_{\mathrm{max}}^2}$$
(26)
This result coincides exactly with the semi-classical estimate given in Ref. . The validity of Eq. (26) can be checked numerically. For instance, for $`x_{max}=7`$ and $`\omega =1`$, we get $`ฯต_0=4.14\times 10^{21}`$ from Eq. (26) which agrees well with $`4.10\times 10^{21}`$ obtained numerically.
In Fig. 4, we have plotted the asymptotic form Eq. (26) but also the full integral form
$$ฯต_0(\sqrt{\omega }x_{max})=1/(2\omega _0^{x_{max}}๐y\mathrm{e}^{\omega y^2}_0^y๐z\mathrm{e}^{\omega z^2}),$$
(27)
One notices that the integral form stays accurate to much lower values of $`x_{max}`$, even at very low values of $`x_{max}`$, where it has no reason to be accurate, it gives a reasonable answer. This observation will be discussed further in Sec. 7.
### 4.2 Anharmonic corrections
Having solved the problem for $`\lambda =0`$ at first order in $`ฯต_n`$, we can use the usual methods of perturbation theory. $`E_0^{(1)}`$ can be calculated by taking the average of $`x^4`$ with the corrected ground state wave function constructed above:
$$E_0^{(1)}=\omega ^2N^1_0^{x_{max}}๐x|\mathrm{\Psi }_0^{(0)}(x)|^2x^4.$$
(28)
with the normalization constant
$$N=_0^{x_{max}}๐x|\mathrm{\Psi }_0^{(0)}(x)|^2.$$
(29)
Proceeding as for Eq. (26), we obtain in leading order
$$E^{(1)}(3/4)\pi ^{1/2}\omega ^{5/2}x_{max}^5\mathrm{e}^{\omega x_{max}^2}$$
(30)
Note that in this calculation, we have replaced $`G_0(x)`$ in the integral by its asymptotic form ($`x^1\mathrm{e}^{\omega x^2}`$ see Eqs. (24 -26)). The exponentials exactly cancel inside the integral and we are left with the integration of $`x^4/x`$. This is justified by the fact that most of the contributions come from the large $`x`$ region. The validity of the Eq. (30) can be checked numerically. For instance, for $`x_{max}=7`$ and $`\omega =1`$, Eq. (30) gives a correction $`4.97\times 10^{18}`$ while numerically we obtain $`5.01\times 10^{18}`$.
The good accuracy of Eq. (30) is shown in Fig. 5. In this figure, we also show the full integral formula of Eq. (28), where the integral was done numerically using the square of the order $`ฯต`$ approximation for the ground state wave function, without discarding order $`ฯต^2`$ terms in the square. Again the integral formula stays valid at unexpectedly low values of $`x_{max}`$. This will be explained in Sec. 7.
The derivation of Eq. (30) is quite simple: $`ฯต_0`$ give a contribution proportional to $`x_{max}e^{\omega x_{\mathrm{max}}^2}`$ and the $`x`$-integral gives a factor $`x_{max}^4`$. We conjecture that at leading order, a similar situation is encountered for higher order terms, and that
$$1R_k(\sqrt{\omega }x_{\mathrm{max}})x_{\mathrm{max}}^{4k+1}e^{\omega x_{\mathrm{max}}^2}.$$
(31)
For this purpose we have studied numerically the quantity
$$Q_k(\sqrt{\omega }x_{max})e^{+\omega x_{\mathrm{max}}^2}(1R_k(\sqrt{\omega }x_{\mathrm{max}})),$$
(32)
which according to the conjecture should scale like $`x_{max}^{4k+1}`$. In Fig. 6, we have set $`\omega =1`$ and displayed $`\mathrm{ln}(Q_k(x_{max})`$ versus $`\mathrm{ln}(x_{max})`$. The solid lines represent the linear functions $`A_k+(4k+1)\mathrm{ln}(x_{max})`$. The constants $`A_k`$ have been fitted using the last ten data points. Their numerical values are $`A_0=0.802`$ which compares well with the prediction of Eq. (26) $`\mathrm{ln}(4\pi ^{1/2})0.814`$, and $`A_1=0.276`$ which compares well with the prediction of Eq. (30) $`\mathrm{ln}(4\pi ^{1/2}/3)0.285`$. The linear behavior of the higher orders supports the conjecture.
## 5 The small $`x_{\mathrm{max}}`$ limit
When the field cutoff $`x_{\mathrm{max}}0`$, the potential term in the harmonic oscillator is much smaller then the kinetic term, and in first approximation we can use the free particle in a box of length $`2x_{\mathrm{max}}`$. The energy levels diverge as $`x_{max}^2`$ and it is convenient to introduce the following rescalings
$$\stackrel{~}{H}=x_{max}^2H$$
$$\stackrel{~}{x}=x/x_{max}$$
$$\stackrel{~}{\omega }=\omega x_{max}^2$$
$$\stackrel{~}{\lambda }=\lambda x_{max}^6$$
We then obtain a new hamiltonian
$$\stackrel{~}{H}=\frac{1}{2}(\frac{d}{d\stackrel{~}{x}})^2+\stackrel{~}{V},$$
(33)
with
$$\stackrel{~}{V}(\stackrel{~}{x})=\{\begin{array}{ccc}\frac{1}{2}\stackrel{~}{\omega }^2\stackrel{~}{x}^2+\stackrel{~}{\lambda }\stackrel{~}{x}^4& \mathrm{if}& |\stackrel{~}{x}|<1\\ \mathrm{}& \mathrm{if}& |\stackrel{~}{x}|1\end{array}$$
(34)
With these rescalings, $`x_{max}`$ has disappeared from the problem and we can now expand $`\stackrel{~}{E}`$ as a double series in $`\stackrel{~}{\lambda }`$ and $`\stackrel{~}{\omega }`$. This can be done using the usual methods of perturbation theory since we can solve the problem exactly when $`\stackrel{~}{\lambda }=\stackrel{~}{\omega }=0`$. The perturbative series becomes
$$\stackrel{~}{E}_n=\underset{k=0,l=0}{\overset{\mathrm{}}{}}E_n^{(k,l)}\stackrel{~}{\lambda }^k\stackrel{~}{\omega }^{2l},$$
(35)
with dimensionless coefficients $`E_n^{(k,l)}`$. Scaling back to the original problem, we obtain
$$E_n=\underset{k=0,l=0}{\overset{\mathrm{}}{}}E_n^{(k,l)}\lambda ^k\omega ^{2l}x_{max}^{6k+4l2}$$
(36)
Comparing with Eq. (9), we conclude that with this expansion
$$E_n^{(k)}(\omega x_{max}^2)=\underset{l=0}{\overset{\mathrm{}}{}}E_n^{(k,l)}\times (\omega x_{max}^2)^{2l+3k1}.$$
(37)
Obviously, this implies that the asymptotic behavior when $`x_{max}0`$, for the ratios displayed in Fig. 1, is
$$R_k(\sqrt{\omega }x_{max})x_{max}^{6k2}$$
(38)
We checked numerically that this asymptotic behavior is correct. Using the data of Fig. 1, one can indeed determine the constant of proportionality $`E_0^{(l,0)}`$ calculated independently below with more than two significant digits.
For the ground state,
$`E_0^{(0,0)}`$ $`=`$ $`{\displaystyle \frac{\pi ^2}{8}}`$
$`E_0^{(1,0)}`$ $`=`$ $`{\displaystyle \frac{1}{5}}{\displaystyle \frac{4}{\pi ^2}}+{\displaystyle \frac{24}{\pi ^4}}`$
$`E_0^{(0,1)}`$ $`=`$ $`{\displaystyle \frac{1}{6}}{\displaystyle \frac{1}{\pi ^2}}`$ (39)
Higher orders require infinite sums. In practice, since we see from Eq. (34) that the calculations can be performed with $`x_{max}=1`$, it is easier to proceed numerically as in Sec. 2 but with a double series. The numerical results are given in Table 1 and displayed in Fig. 7.
The approximately linear behavior of $`ln(|E_0^{(k,l)}|)`$ in $`l`$ at fixed $`k`$ shown in Fig. 7 suggests that the power series in $`\omega ^2x_{max}^4`$, for the coefficients of $`(\lambda /\omega ^3)^k`$ have a finite radius of convergence (see Eq. (37)). As the slope vary from -4.7 for $`k=0`$ to -3.3 for $`k=5`$, we expect that the series converge for $`|x_{max}\sqrt{\omega }|<\mathrm{e}^{4.7/4}3.3`$ for $`k=0`$ to $`\mathrm{e}^{3.3/4}2.3`$ for $`k=5`$. Fig. 8 shows a steady decrease of the radius of convergence when $`k`$ increases. At the same time, the value of $`x_{max}`$ where the transition between the small and large cutoff limit occurs, denoted $`x_0(k)`$ increases with $`k`$. This quantity is defined more precisely in Sec. 7. Consequently, it seems clear that as the order increases, a gap opens between the range of validity of the small $`x_{max}`$ expansion and the crossover region.
## 6 Rigorous lower bounds ont radius of convergence
In Secs. 3 and 5, we have discussed the radius of convergence of series from a numerical point of view. We now return to this question with more analytical considerations. First of all, it is possible to get an order magnitude estimate for the radius of convergence by using a simple argument. If we look at a typical term of the perturbative series at order $`\lambda ^k`$, it is the product of $`k`$ matrix elements for the perturbation potential $`x^4`$ divided by $`k1`$ differences of energy level. The matrix elements of $`x^4`$ are bounded in magnitude by $`x_{max}^4`$. The differences of energy depend on $`x_{max}`$. If $`x_{max}`$ is very small, in other words if $`x_{max}^2\omega ^2<<\omega `$ in units where the mechanical mass and $`\mathrm{}`$ are set to 1, the harmonic term is a perturbation and the energy differences are proportional to $`x_{max}^2`$. For instance, the difference between the ground state and the first excited state will be approximately $`3\pi ^2/(8x_{max}^2)`$. In the opposite limit of large $`x_{max}`$, the level are approximately equidistant with a separation $`\omega `$. From this we expect that the perturbative series converges for $`|\lambda |<\lambda _c(x_{max},\omega )`$ with
$$\lambda _c(x_{max},\omega )\{\begin{array}{cc}x_{max}^6\hfill & \mathrm{if}x_{max}<<\omega ^{1/2}\hfill \\ \omega x_{max}^4\hfill & \mathrm{if}x_{max}>>\omega ^{1/2}\hfill \end{array}$$
(40)
This argument does not take into account aspects such as the number of terms present or their signs, however, a rigorous lower bound on the radius of convergence can be obtained by simply using one half of the difference between the ground state and the first excited state of the unperturbed system in the above argument. This is formulated precisely in Theorem XII.11 in Ref. and proved in Ref. . The final result in our case is
$$\lambda _c(x_{max},\omega )>(1/2)(E_1^{(0)}E_0^{(0)})/x_{max}^4$$
(41)
Note that Theorem XII.11 is more than what we need here because the perturbation is a bounded operator and we do not need to compare its norm with the norm of $`H_0`$.
In the case $`\omega =0`$, dimensional analysis dictates $`\lambda _c(x_{max},\omega =0)=Cx_{max}^6`$ for some positive $`C`$, and Eq. (41) implies
$$\lambda _c(x_{max},\omega =0)>(3\pi ^2/16)\times x_{max}^6$$
(42)
This bound can be compared with numerical estimates obtained from the series with coefficients $`E_0^{(k,0)}`$ defined in Sec. 5 and which only requires one calculation with $`x_{max}=1`$ instead of the two step procedure of Sec. 3. Using the coefficients for $`k=5,\mathrm{}10`$, we obtain
$$\lambda _c(x_{max},\omega =0)53\times x_{max}^6,$$
(43)
which is quite close to result obtained with $`\omega =1`$ at low $`x_{max}`$ in Sec. 5. The rigorous bound on $`C`$, $`3\pi ^2/161.85`$ is about 30 times smaller than the empirical value 53. In summary, the lower bound is compatible with our numerical estimate, but is not sharp.
The same argument can be applied for the calculation of $`E_0^{(0)}`$ as a power series in $`(1/2)\omega ^2`$ for the perturbation $`(1/2)\omega ^2x^2`$. It can be tested by studying the coefficients $`E^{(0,l)}`$ as we did in Sec. 5. The rigorous bound becomes $`(1/2)\omega _c^2>(1/2)(3\pi ^2/8)/x_{max}^4`$, which implies
$$\omega _c>1.92\times x_{max}^2$$
(44)
From Sec. 5, we have
$$\omega _c(3.3)^2\times x_{max}^2,$$
(45)
and again, the lower bound is compatible with the numerical estimate, but is not an accurate estimator of the actual radius of convergence.
## 7 The crossover region: empirical data collapse
In this section, we discuss the question of the interpolation between the large and small $`x_{max}`$ approximation. In Sec. 4, we have already observed that by naively using integral formulas derived in the large $`x_{max}`$ approximation for $`E^{(0)}`$ and $`E^{(1)}`$, it was possible to obtain decent approximations at low $`x_{max}`$. We will first explain how this works and then discuss the higher orders terms.
In the (mathematical) limit of very small $`x_{max}`$, Eq. (23) implies that $`G_0(x)\omega x^2`$. Consequently, in this limit, where we do not expect the approximation to be accurate, the correction becomes dominant and we have
$$E_0^{(0)}(\sqrt{\omega }x_{max})1/(\omega x_{max}^2).$$
(46)
Despite the lack of justification for this formula, it is not far from the accurate answer derived in Sec. 5:
$$E_0^{(0)}(\sqrt{\omega }x_{max})\pi ^2/(8\omega x_{max}^2)1.234/(\omega x_{max}^2)$$
(47)
Similarly for $`E^{(1)}`$, we can study Eq. (28) in the mathematical limit of small $`x_{max}`$. According to Sec. 4, $`\mathrm{\Psi }_0^{(0)}(x)(1(x/x_{max})^2)`$ and elementary integration yields
$$E_0^{(1)}(\sqrt{\omega }x_{max})(1/21)\times (\omega x_{max}^2)^20.0476(\times \omega x_{max}^2)^2$$
(48)
Again, this is close to the accurate answer from Eq. (5):
$$E_0^{(1)}(\sqrt{\omega }x_{max})0.0411\times (\omega x_{max}^2)^2$$
(49)
For higher orders, we know from Fig. 8 that the small $`x_{max}`$ approximation is not convergent in the crossover region. However, the shape similarities observed in Fig. 1 suggest to parametrize $`R_k`$ in terms of a single function $`U`$, that can be shifted by a $`k`$-dependent quantity that we denote $`x_0(k)`$. We have chosen $`x_0(k)`$ as the value of $`x_{max}`$ for which the second derivative of $`R_k`$ vanishes. The numerical values are shown in Table 2 and Fig. 8.
Empirically, it can be fitted quite well with
$$x_0(k)0.87+1.13\sqrt{k}.$$
(50)
In Fig. 9, we show that the possibility of having
$$R_k(x_{max})U(x_{max}x_0(k)),$$
(51)
is reasonably well satisfied for $`k2`$. In other words, with suitable translations, the $`R_k`$ approximately โcollapseโ.
## 8 Conclusions
We have shown that the calculation of the perturbative series for the anharmonic oscillator with a field cutoff $`x_{max}`$ can be performed reliably using numerical methods, and approximate analytical methods in the limit of small and large $`x_{max}`$. For the coefficients of order 0 and 1 in $`\lambda `$, the integral formulas derived in the large $`x_{max}`$ limit, produce the correct leading power dependence with coefficients close to the correct ones, in the opposite limit (small $`x_{max}`$) where they are not expected to be accurate. For higher orders, it is possible to approximately collapse the crossover region of the various order by an appropriate translation in $`x_{max}`$. The collapse is not perfect and it is necessary to resolve more accurately the details of Fig. 9 in order to obtain accurate formulas for the $`R_k`$.
We expect that it is possible to extend the small and large cutoff techniques in the case of higher dimensional field theory, however the only numerical method that can we envision for the crossover region is the Monte Carlo method . In view of this, it is essential to reach a proper understanding of the interpolation.
This research was supported in part by the Department of Energy under Contract No. FG02-91ER40664 and in part by the National Science Foundation under Grant No. PHY99-07949. This work was completed while Y. M. was visiting the Kavli Institute for Theoretical Physics. Y. M. thanks the organizers and participants of the workshop โModern Challenges for Lattice Field Theoryโ for conversations. We thank J. Cook for checking series expansions in Sec. 4 and B. Kessler for providing independent checks of the numerical estimates of the radius of convergence.
## References
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# 1 Introduction
## 1 Introduction
We suppose that we are given probability spaces $`(S,๐ฎ,\mu )`$ and $`(E,,\nu )`$ and a mapping
(1.1)
$$\varphi :S\times S\times ES$$
under which the image of the product measure $`\mu \mu \nu `$ is $`\mu `$. In general $`\varphi :S^n\times ES`$ may be considered, even with $`n`$ being infinity, but here we take $`n=2`$ for simplicity and we further assume that $`\varphi `$ is symmetric in its first two arguments. Typically in applications $`\nu `$ and $`\varphi `$ are given and the measure $`\mu `$ is to be determined. $`\mu `$ is then said to be a solution to a recursive distributional equation. Existence and uniqueness of $`\mu `$ is well studied in a number of important cases; the survey paper by Aldous and Bandyopdhyay, , describes many examples. Here our concerns are somewhat different. Associated with this setup is a recursive tree process constructed as follows. Let $`\mathrm{\Gamma }_{\mathrm{}}=_{k0}\{0,1\}^k`$ be the set of vertices of the infinite rooted binary tree; in particular $`\mathrm{}\mathrm{\Gamma }_{\mathrm{}}`$ denotes the root. The set of vertices on level $`n`$ of the tree, $`\{0,1\}^n`$, is denoted by $`G_n`$, and $`\mathrm{\Gamma }_n=_{k=0}^n\{0,1\}^k`$ denotes the set of vertices belonging to levels up to and including $`n`$. A vertex $`u=(u_1,u_2,\mathrm{},u_n)`$ has two daughters: vertices $`u0=(u_1,\mathrm{},u_n,0)`$ and $`u1=(u_1,\mathrm{},u_n,1)`$. Let $`\left(\xi _u;u\mathrm{\Gamma }_{\mathrm{}}\right)`$ and $`\left(ฯต_u;u\mathrm{\Gamma }_{\mathrm{}}\right)`$ be the co-ordinate maps on $`S^\mathrm{\Gamma }_{\mathrm{}}\times E^\mathrm{\Gamma }_{\mathrm{}}`$. Define
(1.2)
$$\mathrm{\Omega }=\{\omega S^\mathrm{\Gamma }_{\mathrm{}}\times E^\mathrm{\Gamma }_{\mathrm{}}:\xi _u(\omega )=\varphi (\xi _{u0}(\omega ),\xi _{u1}(\omega ),ฯต_u(\omega ))\text{ for each }u\mathrm{\Gamma }_{\mathrm{}}\},$$
and let $``$ denote the restriction of the product $`\sigma `$-algebra to $`\mathrm{\Omega }`$. There is a unique probability measure $`m`$ on $`(\mathrm{\Omega },)`$ under which for each $`n`$, the joint law of $`\left(\xi _u;uG_n\right)`$ and $`\left(ฯต_u;u\mathrm{\Gamma }_{n1}\right)`$ is $`^{G_n}\mu ^{\mathrm{\Gamma }_{n1}}\nu `$. We think of independent random variables $`\left(ฯต_u;u\mathrm{\Gamma }_{\mathrm{}}\right)`$, sometimes referred to as the innovations, as the input that drives the system, while $`\left(\xi _u;u\mathrm{\Gamma }_{\mathrm{}}\right)`$ is thought of as the response. Aldous and Bandyopadhyay draw attention to investigating when the recursive tree process is endogenous, that is to say every random variable on $`(\mathrm{\Omega },)`$ is $`m`$-almost surely equal to a function of the innovations process $`\left(ฯต_u;u\mathrm{\Gamma }_{\mathrm{}}\right)`$ alone. Loosely speaking this means there is no additional randomness in the system โlocatedโ at the boundary of the infinite tree. However care must be taken in interpreting this: the tail $`\sigma `$-algebra
(1.3)
$$\underset{n}{}\sigma \left(\xi _u,ฯต_u;u\mathrm{\Gamma }_{\mathrm{}}\mathrm{\Gamma }_n\right)$$
is typically empty even if endogeny does not hold! A strong parallel may be drawn with certain stochastic differential equations which admit weak but not strong solutions, see in particular Tsirelsonโs example discussed in Section V.18 of . In non-endogenous cases it is natural to try to give some explicit description of the additional randomness, however this does not seem possible in any generality. One case that is amenable to such an analysis is the linear โsmoothing transformationโ which has been extensively studied. Various representations theorems obtained by Durrett and Liggett, , Liu , Caliebe and Rรถsler , and others reveal a structure in which the innovations are augmented with Gaussian or Poisson noise living on the boundary of the tree. Another case whose structure can similarly be understood is studied by Biggins .
Let $`u_0=\mathrm{},u_1,u_2,\mathrm{},u_n,\mathrm{}`$ be an infinite sequence of vertices, $`u_{n+1}`$ being a daughter of $`u_n`$ for each $`n`$. For $`n0`$ let $`\xi _n=\xi _{u_n}`$. The law of the sequence $`\left(\xi _n;n0\right)`$, which by the symmetry of $`\varphi `$ does not depend on the choice of sequence of vertices, is easily seen to be that of a stationary Markov chain (indexed by negative time) with transition kernel $`P`$ defined on $`S\times ๐ฎ`$ given by
(1.4)
$$P(x_0,A)=_S_E\mathrm{๐}\left(\varphi (x_0,x_1,z)A\right)\nu (dz)\mu (dx_1).$$
In this paper we give a criterion for endogeny is in terms of the corresponding โtwo point motionโ, that is to say a Markov chain on $`S^2`$ with transition kernel $`P^{(2)}`$ on $`S^2\times ๐ฎ^2`$ given by
(1.5)
$$P^{(2)}((x_0,x_0^{}),A\times A^{})=_S_E\mathrm{๐}\left(\varphi (x_0,x_1,z)A,\varphi (x_0^{},x_1,z)A^{}\right)\nu (dz)\mu (dx_1).$$
Let $`S^{}`$ be the diagonal of $`S^2`$ which is an absorbing set for $`P^{(2)}`$. Let $`P^{()}`$ be the restriction of $`P^{(2)}`$ to $`S^2S^{}`$.
We now assume that $`S`$ is finite.
In this case we can identify $`P^{()}`$ with a non-negative square matrix to which we can apply Perron-Frobenius theory. Let $`\rho `$ be the largest eigenvalue of $`P^{()}`$.
###### Theorem 1.
If $`2\rho <1`$ then the tree process is endogenous, whereas if $`2\rho >1`$ then the process is not endogenous. In the critical case $`2\rho =1`$, if the additional conditions (4.1) and (4.2) hold, then the tree process is endogenous.
Aldous and Bandyopadhyay have given a different necessary and sufficient condition for endogeny, the above result being related to their condition by linearization. We now elucidate this connection. Associated with $`\varphi :S\times S\times ES`$ is the โtwo-point mapโ $`\varphi ^{(2)}:S^2\times S^2\times ES^2`$ given by
(1.6)
$$\varphi ^{(2)}((x_0,x_0^{}),(x_1,x_1^{}),z)=(\varphi (x_0,x_1,z),\varphi (x_0^{},x_1^{},z)).$$
Now we can define $`T^{(2)}:(S^2)\times (S^2)(S^2)`$ by setting $`T^{(2)}(\lambda _0,\lambda _1)`$ equal to the image of the measure $`\lambda _0\lambda _1\nu `$ under the map $`\varphi ^{(2)}`$. Recall that $`\mu `$ denotes our given measure on $`S`$ that is invariant for the recursive distributional equation corresponding to $`\varphi `$ and $`\nu `$. The measure $`\mu ^{}`$ on $`S^2`$, carried by the diagonal and having marginals $`\mu `$, is a fixed point of $`T^{(2)}`$ in that $`T^{(2)}(\mu ^{},\mu ^{})=\mu ^{}`$. According to Theorem 11 of , under some minor technical condition, $`\mu ^{}`$ being the only probability measure on $`S^2`$ that is a fixed point of $`T^{(2)}`$, and whose marginals are both equal to $`\mu `$, is a necessary and sufficient condition for endogeny. Bandyopadhyay, , has used this criterion to show the endogeny of certain important examples that arise from applications. The map $`T^{(2)}`$ is bilinear and symmetric so that we have for real $`\delta `$ and $`\lambda (S^2)`$
(1.7)
$$T^{(2)}(\mu ^{}+\delta \lambda ,\mu ^{}+\delta \lambda )=\mu ^{}+2\delta T^{(2)}(\lambda ,\mu ^{})+\delta ^2T^{(2)}(\lambda ,\lambda ),$$
which shows that $`2T^{(2)}(,\mu ^{}):(S^2)(S^2)`$ is the derivative of $`\lambda T^{(2)}(\lambda ,\lambda )`$ at the fixed point $`\mu ^{}`$. It is straightforward to verify that
(1.8)
$$\lambda P^{(2)}=T^{(2)}(\lambda ,\mu ^{}).$$
Thus the criterion given in Theorem 1 can be viewed in terms of the stability of the fixed point $`\mu ^{}`$ of $`T^{(2)}`$, and it is quite natural that the existence of other fixed points of $`T^{(2)}`$ should be related to this.
Let $`=L^2(\mathrm{\Omega },,m)`$, and $`๐ฆ`$ be the subspace of (equivalence classes of) random variables measurable with respect to the innovations process $`\left(ฯต_u;u\mathrm{\Gamma }_{\mathrm{}}\right)`$. Endogeny means $`=๐ฆ`$. We also have subspaces $`๐ฆ_n`$ and $`_n`$ containing random variables measurable with respect to $`\left(ฯต_u;u\mathrm{\Gamma }_{n1}\right)`$ and with respect to $`\left(ฯต_u;u\mathrm{\Gamma }_{n1}\right)`$ together with $`\left(\xi _u;uG_n\right)`$ respectively.
In the following three sections we treat the three cases: subcritical $`2\rho <1`$; supercritical $`2\rho >1`$; and critical $`2\rho =1`$ separately. The subcritical case is treated by using an operator version of Markovโs inequality to bound from above the quantity $`f^2P_{๐ฆ_n}f^2`$ where $`P_{๐ฆ_n}f`$ is the orthogonal projection of $`f`$ onto $`๐ฆ_n`$. The supercritical case is treated by using the left eigenvector of $`P^{()}`$ corresponding to $`\rho `$ to construct a consistent family of quadratic forms on $``$, and hence an operator $`Q_{\mathrm{}}`$. This operator is the generator of dynamics that fix the innovations while perturbing the additional randomness at the boundary. The existence of such dynamics precludes endogeny.
In the final section of the paper we show that the dynamics associated with $`Q_{\mathrm{}}`$ naturally arise through a passage to the limit. Consider for each $`n1`$ the finite configuration space
(1.9)
$$\mathrm{\Omega }_n=\{\omega S^{\mathrm{\Gamma }_n}\times E^{\mathrm{\Gamma }_{n1}}:\xi _u(\omega )=\varphi (\xi _{u0}(\omega ),\xi _{u1}(\omega ),ฯต_u(\omega ))\text{ for each }u\mathrm{\Gamma }_{n1}\},$$
which we equip with the measure $`m_n`$ under which the joint law of $`\left(\xi _u;uG_n\right)`$ and $`\left(ฯต_u;u\mathrm{\Gamma }_{n1}\right)`$ is $`^{G_n}\mu ^{\mathrm{\Gamma }_{n1}}\nu `$. Consider the dynamics on $`\mathrm{\Omega }_n`$ with invariant measure $`m_n`$ under which the co-ordinates $`\xi _u`$ for $`uG_n`$ are independently refreshed at rate $`1`$, whilst $`\xi _v`$ for $`v\mathrm{\Gamma }_{n1}`$ are determined by application of the map $`\varphi `$ using the innovations $`(ฯต_u;u\mathrm{\Gamma }_{n1})`$ which are held fixed for all time. Let $`A_n`$ denote the corresponding generator acting on $`L^2(\mathrm{\Omega }_n,m_n)`$. We show that at the level of generators, these dynamics, when slowed down by a factor of $`(2\rho )^n`$, converge to those associated with $`Q_{\mathrm{}}`$. Let $`P__n:_n`$ be the orthogonal projection onto $`_n`$ which we identify with $`L^2(\mathrm{\Omega }_n,m_n)`$.
###### Theorem 2.
Suppose that $`P^{()}`$ is primitive, and that $`2\rho >1`$, then
$$(2\rho )^nA_nP__nQ_{\mathrm{}},$$
in the strong resolvent sense as $`n`$ tends to infinity.
There is strong similarity between the methods used in this paper and those used to study multitype branching processes. To make this connection fix $`f_0`$ that satisfies $`f^2=1`$ and consider the probability distribution distributions $`\mu _f^{(n)}`$ on the non-negative integers defined by
(1.10)
$$\mu _f^{(n)}(k)=P_{E_k}f^2\text{ for }k0,$$
where $`P_{E_k}`$ is the projection operator associated with the eigenspace $`E_k`$ of the operator $`A_n`$ corresponding to the eigenvalue $`k`$. See below at (2.3) for further information on this. We use arguments that are based on treating this distribution as if it were that of the number of particles alive in generation $`n`$ of a branching process. In particular Theorem 1 corresponds exactly to the fact the certainty or otherwise of eventual extinction for a branching process depends upon whether the Malthusian parameter of the process is greater than one. The arguments used are also closely related to the spectral methods employed by Tsirelson, , to study continuous products of probability spaces. In fact Tsirelson and Vershik introduced certain recursive tree processes in as examples in the theory of continuous products.
## 2 The subcritical case: endogeny
Recall the probability measure $`m`$ on $`(\mathrm{\Omega },)`$ is characterized by the fact that the joint law of $`\left(\xi _u;uG_n\right)`$ and $`\left(ฯต_u;u\mathrm{\Gamma }_{n1}\right)`$ is $`^{G_n}\mu ^{\mathrm{\Gamma }_{n1}}\nu `$. Thus there are isomorphisms between Hilbert spaces:
(2.1)
$$\begin{array}{c}_nL^2(S^{G_n}\times E^{\mathrm{\Gamma }_{n1}},๐ฎ^{G_n}\times ^{\mathrm{\Gamma }_{n1}},^{G_n}\mu ^{\mathrm{\Gamma }_{n1}}\nu )\hfill \\ \hfill ^{G_n}L^2(S,๐ฎ,\mu )^{\mathrm{\Gamma }_{n1}}L^2(E,,\nu ).\end{array}$$
$`L^2(S,๐ฎ,\mu )`$ is the direct sum of the one-dimensional subspace of constants together with its orthogonal complement to be denoted by $`L_0^2(S,๐ฎ,\mu )`$. Decomposing each copy of $`L^2(S,๐ฎ,\mu )`$ appearing on the righthandside of (2.1) in this manner we obtain
(2.2)
$$_n=\underset{SG_n}{}_S,$$
where for each subset $`SG_n`$ of vertices on level $`n`$ of the tree the corresponding subspace $`_S`$ of the Hilbert space $`_n`$ is generated by vectors of the form $`_{uG_n}f_u_{v\mathrm{\Gamma }_{n1}}g_v`$ with $`f_uL_0^2(S,๐ฎ,\mu )`$ for $`uS`$, $`f_u`$ being a constant vector for $`uS`$, and $`g_vL^2(E,,\nu )`$ being unrestricted.
Given a linear operator $`L`$ acting on $`L^2(S,๐ฎ,\mu )`$ and a vertex $`uG_n`$ we may consider โ$`L`$ applied at $`u`$โ. More precisely we define an operator $`L^{(u)}`$ acting on $`_n`$ as being unitary equivalent, via the isomorphism (2.1), to the tensor product of $`L`$ acting on the copy of $`L^2(S,๐ฎ,\mu )`$ corresponding to $`u`$ together with the identity on all other factors. We consider the case that $`L`$ is given by $`P_\mathrm{๐}^{}`$ which is the orthogonal projection onto the subspace $`L_0^2(S,๐ฎ,\mu )`$. We define an operator on $`_n`$ via
(2.3)
$$A_n=\underset{uG_n}{}P_\mathrm{๐}^{}^{(u)}.$$
By considering its action on the generating vectors for each subspace $`_S`$ we find that the eigenvalues of $`A_n`$ are $`0,1,2,\mathrm{},2^n`$, with the eigenvalue $`k`$ having corresponding eigenspace $`_{|S|=k}_S`$. Here $`|S|`$ denotes the number of vertices belonging to $`S`$. On the other hand the operator $`IP_{๐ฆ_n}`$, which commutes with $`A_n`$, has eigenvalues $`0`$ and $`1`$ with corresponding eigenspaces $`๐ฆ_n=_{\mathrm{}}`$ and $`_S\mathrm{}_S`$. By decomposing $`f_n`$ according to the subspaces $`_S`$ we deduce the following inequality.
###### Proposition 3.
For every $`f_n`$,
$$0(IP_{๐ฆ_n})f^2(f,A_nf).$$
###### Proof of Theorem 1: Subcritical case..
Because of the recursive structure endogeny is is equivalent to $`_0๐ฆ`$. So fix $`f_0`$, which, in a slight abuse of notation, we also treat as an element of $`L^2(S,๐ฎ,\mu )`$. To prove $`f`$ belongs to $`๐ฆ`$ it is enough, by virtue of Proposition 3, to prove that
$$(f,A_nf)0\text{ as }n\mathrm{}.$$
We may express $`(f,A_nf)`$ using couplings of the tree-indexed process. Fix a vertex $`u_nG_n`$. Let $`\mathrm{\Omega }_n^{}`$ be a copy of the finite configuration space $`\mathrm{\Omega }_n`$ and consider the product space $`\mathrm{\Omega }_n\times \mathrm{\Omega }_n^{}`$ equipped with co-ordinate maps $`(ฯต_u,ฯต_u^{};u\mathrm{\Gamma }_{n1})`$ and $`(\xi _u,\xi _u^{};u\mathrm{\Gamma }_n)`$. Let the probability measure $`\stackrel{~}{m}_n`$ on this product space have both marginals equal to $`m_n`$ and be such that
$`\stackrel{~}{m}_n`$ is supported on the set where $`ฯต_u=ฯต_u^{}`$ for all $`u\mathrm{\Gamma }_{n1}`$ and $`\xi _u=\xi _u^{}`$ for all $`uG_n`$ except $`u_n`$;
$`(ฯต_u;u\mathrm{\Gamma }_{n1})`$, $`(\xi _u;u\mathrm{\Gamma }_n)`$ and $`\xi _{u_n}^{}`$ are independent under $`\stackrel{~}{m}_n`$.
It is easily verified that
$$(f,P_\mathrm{๐}^{}^{(u_n)}f)=\frac{1}{2}_{\mathrm{\Omega }_n\times \mathrm{\Omega }_n^{}}(f\xi _{\mathrm{}}f\xi _{\mathrm{}}^{})^2๐\stackrel{~}{m}_n=(\mu \mu )P_n^{(2)}g,$$
where $`P_n^{(2)}`$ is the $`n`$-step transition matrix for $`P^{(2)}`$ and $`g`$ is the function $`g(x,x^{})=\frac{1}{2}(f(x)f(x^{}))^2`$ on $`S^2`$. Summing over the possible choices of $`u_n`$ we obtain
$$(f,A_nf)=2^n(\mu \mu )P_n^{(2)}g.$$
Since $`g`$ is zero on the diagonal of $`S^2`$ this quantity can also be expressed using $`P^{()}`$, and then, since the spectral radius of $`P^{()}`$ is less than $`\frac{1}{2}`$ by hypothesis, we obtain the desired convergence to zero as $`n`$ tends to infinity. โ
## 3 Dynamics and non-endogeny
Let $`L`$ be the operator on $`L^2(S,๐ฎ,\mu )`$ associated with a matrix $`\left(L(x,x^{});x,x^{}S\right)`$,
$$Lf(x)=\underset{x^{}S}{}L(x,x^{})f(x^{}).$$
We use the transition probabilities $`P^{(2)}`$ to determine a new operator $`๐ซL`$ also acting on $`L^2(S,๐ฎ,\mu )`$, which is defined via its associated matrix via,
(3.1)
$$(\mu L)P^{(2)}=\mu (๐ซL).$$
Here $`(\mu L)(x,x^{})=\mu (x)L(x,x^{})`$ is treated as a row vector on $`S^2`$ and $`P^{(2)}`$ acts by matrix multiplication on the right. In a similar way we may define a quadratic superoperator $`๐ฌ`$ by using the mapping $`T^{(2)}`$,
(3.2)
$$T^{(2)}(\mu L,\mu L)=\mu (๐ฌL).$$
Both superoperators $`๐ซ`$ and $`๐ฌ`$ arise when considering the isometric embedding of Hilbert spaces: $`L^2(S,๐ฎ,\mu )L^2(S^2\times E,๐ฎ^2\times ,^2\mu \nu )`$ given by $`ff\varphi `$. Given $`L`$ acting on on $`L^2(S,๐ฎ,\mu )`$, it is easily verified that the new operator $`๐ซL`$ satisfies
(3.3)
$$(f\varphi ,(LII)g\varphi )_{L^2(S^2\times E,๐ฎ^2\times ,^2\mu \nu )}=(f,(๐ซL)g)_{L^2(S,๐ฎ,\mu )},$$
for all $`f,gL^2(S,๐ฎ,\mu )`$. Similarly $`L`$ and $`๐ฌL`$ satisfy
(3.4)
$$(f\varphi ,(LLI)g\varphi )_{L^2(S^2\times E,๐ฎ^2\times ,^2\mu \nu )}=(f,(๐ฌL)g)_{L^2(S,๐ฎ,\mu )},$$
for all $`f,gL^2(S,๐ฎ,\mu )`$. It is a consequence of the recursive structure that these two relations extend as is recorded in the following proposition whose proof we omit.
###### Proposition 4.
If $`vG_{n+1}`$ is a daughter of some $`uG_n`$ then, for all $`f,g_n`$,
$$(f,L^{(v)}g)=(f,(๐ซL)^{(u)}g).$$
If $`v_1,v_2G_{n+1}`$ are daughters of distinct $`u_1,u_2G_n`$, then for $`f,g_n`$,
$$(f,L^{(v_1)}L^{(v_2)}g)=(f,(๐ซL)^{(u_1)}(๐ซL)^{(u_2)}g).$$
If $`v_1,v_2G_{n+1}`$ are the two daughters of some $`uG_n`$, then for $`f,g_n`$,
$$(f,L^{(v_1)}L^{(v_2)}g)=(f,(๐ฌL)^{(u)}g).$$
Corresponding to the principal eigenvalue $`\rho `$ of $`P^{()}`$ is a left eigenvector $`\kappa `$ satisfying $`\kappa P^{()}=\rho \kappa `$. Of course $`\kappa `$ is only determined up to a scalar multiple and we make some arbitrary choice. Considered as a function on $`SS^{}`$, $`\kappa `$ is symmetric since $`P^{()}`$ preserves the space of vectors with this symmetry. It is conceivable that $`P^{()}`$ is not irreducible in which case there may be some further freedom in choosing $`\kappa `$. This does matter so long as we always choose it, as we may, to be symmetric. Next we define a symmetric operator $`Q`$ on $`L^2(S,๐ฎ,\mu )`$ from $`\kappa `$ via
(3.5)
$$(f,Qg)_{L^2(S,๐ฎ,\mu )}=\frac{1}{2}\underset{(x,x^{})S^2S^{}}{}\left(f(x^{})f(x)\right)\left(g(x^{})g(x)\right)\kappa (x,x^{}).$$
$`Q`$ is the generator of an $`S`$-valued Markov process which jumps from $`x`$ to $`x^{}`$ at rate $`Q(x,x^{})=\kappa (x,x^{})/\mu (x)`$. We can always assume that $`\mu (x)>0`$ for all $`xS`$ by deleting part of $`S`$ if necessary.
Define an operator on $`_n`$ via
(3.6)
$$Q_n=(2\rho )^n\underset{uG_n}{}Q^{(u)}.$$
This operator is the generator of a Markov process taking values in the finite configuration space $`\mathrm{\Omega }_n`$. The coordinates $`\xi _u`$, with $`uG_n`$, evolve independently, each a copy of the process generated by $`Q`$ but with their speed altered by the factor $`(2\rho )^n`$. At any instant the $`\xi _u`$ coordinates for $`u\mathrm{\Gamma }_{n1}`$ are determined from the $`\xi _u`$ coordinates with $`uG_n`$ by application of the map $`\varphi `$ with the innovations $`(ฯต_u;u\mathrm{\Gamma }_{n1})`$ fixed for all time.
The significance of the family of generators $`Q_n`$ is that they have a certain consistency property that manifests itself at the level of the corresponding forms. We introduce the forms $`_n`$ defined $`_n`$ via
(3.7)
$$_n(f,g)=(f,Q_ng)\text{ for }f,g_n.$$
###### Lemma 5.
The operator $`Q`$ satisfies $`๐ซQ=\rho Q,`$ and consequently the forms $`_n`$ are consistent in the sense that for any $`mn`$,
$$_m(f,g)=_n(f,g)\text{ for all }f,g_m.$$
###### Proof.
In view of the relation between $`๐ซ`$ and $`P^{(2)}`$ given by (3.1), to prove the first assertion we must verify that
$$(\mu Q)P^{(2)}=\rho (\mu Q).$$
$`(\mu Q)(x,x^{})=\kappa (x,x^{})`$ for $`xx^{}`$ and so the desired equality holds on $`SS^{}`$ since there it becomes $`\kappa P^{()}=\rho \kappa `$. We deduce that the equality must also hold on the diagonal by observing that $`๐ซ`$ preserves the class of operators satisfying $`L1=0`$.
Suppose that $`f,g_m`$ then using what we have just shown together with Proposition 4 we obtain
$$\begin{array}{c}_{m+1}(f,g)=(f,Q_{m+1}g)=(2\rho )^{(m+1)}\underset{vG_{m+1}}{}(f,Q^{(v)}g)\hfill \\ \hfill =2(2\rho )^{(m+1)}\underset{uG_m}{}(f,(๐ซQ)^{(u)}g)=(2\rho )^m\underset{uG_m}{}(f,Q^{(u)}g)=_m(f,g),\end{array}$$
which proves the consistency of the forms.
It follows from the consistency of the $`_n`$ just established that we can define a form $``$ on the dense subspace $`_n_n`$ of $``$ via $`(f,g)=_n(f,g)`$ whenever $`f,g_n`$. However it is not necessarily true that $``$ is closable.
###### Lemma 6.
If $`2\rho >1`$ then for $`f_m`$ for some $`m`$,
$$\underset{nm}{sup}Q_nf<\mathrm{}.$$
###### Proof.
Consider $`f_n`$. Expanding $`Q_{n+1}`$ as a sum and using Proposition 4, plus $`๐ซQ=\rho Q`$, gives
$$\begin{array}{c}(f,Q_{n+1}^2f)=(2\rho )^{2n2}\underset{u,vG_{n+1}}{}(f,Q^{(u)}Q^{(v)}f)\hfill \\ \hfill =2(2\rho )^{2n2}\underset{uG_n}{}(f,(๐ฌQ+๐ซQ^2)^{(u)})f)+4(2\rho )^{2n2}_{\begin{array}{c}u,vG_n\\ uv\end{array}}(f,(๐ซQ)^{(u)}(๐ซQ)^{(v)}f)\\ \hfill =2(2\rho )^{2n2}\underset{uG_n}{}(f,(๐ฌQ+๐ซQ^2)^{(u)}f)(2\rho )^{2n}\underset{uG_n}{}(f,(Q^2)^{(u)})f)+(f,Q_n^2f)\\ \hfill =(2\rho )^{2n}\underset{uG_n}{}(f,\widehat{Q}^{(u)}f)+(f,Q_n^2f),\end{array}$$
where the operator $`\widehat{Q}`$ acting on $`L^2(S,๐ฎ,\mu )`$ is given by
$$\widehat{Q}=2(2\rho )^2\left(๐ฌQ+๐ซQ^2\right)Q^2.$$
But now we compare the operator $`_{uG_n}\widehat{Q}^{(u)}`$ with the number operator $`A_n`$ defined by (2.3). Notice that the constant $`1L^2(S,๐ฎ,\mu )`$ satisfies $`\widehat{Q}1=0`$. Thus each subspace $`_S`$ of $`_n`$ is an invariant subspace for $`_{uG_n}\widehat{Q}^{(u)}`$ whose restriction to $`_S`$ has norm $`|S|\widehat{Q}`$. Thus we have
$$\underset{uG_n}{}(f,\widehat{Q}^{(u)}f)\widehat{Q}(f,A_nf).$$
If we express $`(f,A_nf)`$ in terms of $`P^{()}`$ as we did in the previous section then we find that it is bounded by some constant times $`(2\rho )^n`$, and thus we deduce that, for an appropriate constant $`C`$,
$$(f,Q_{n+1}^2f)(f,Q_n^2,f)+C(2\rho )^n.$$
If $`2\rho >1`$ then the desired conclusion follows. โ
###### Proof of Theorem 1: Supercitical case..
The consistency of the forms $`_n`$ can be expressed in terms of the corresponding operators as $`Q_mf=P__mQ_nf`$ for $`f_m`$ and $`nm`$. Consequently if $`sup_{nm}Q_mf<\mathrm{}`$ then, as $`n`$ tends to infinity, $`Q_nf`$ converges in $``$ to some limit we denote by $`\stackrel{ฬ}{Q}_{\mathrm{}}f`$. We see that $`\stackrel{ฬ}{Q}_{\mathrm{}}`$ is a positive symmetric operator with
$$(f,g)=(f,\stackrel{ฬ}{Q}_{\mathrm{}}g)\text{ for }f,g\underset{n}{}_n,$$
and consequently $``$ is closable \[ see Theorem X.23 of \].
By construction the subspace $`๐ฆ_n`$ lies in the kernel of the operator $`Q_n`$ for each $`n`$. Hence
$$(f,f)=0\text{ for all }f\underset{n}{}๐ฆ_n.$$
Denoting the closure of $``$ by $`\overline{}`$ we deduce that
$$\overline{}(f,f)=0\text{ for all }f๐ฆ.$$
But $`\overline{}(f,f)`$ cannot be identically zero on $``$ since $`Q`$, and hence $`_0`$, was not zero, thus $`๐ฆ`$. โ
## 4 The critical case
The critical case $`2\rho =1`$ is endogenous, provided we impose two additional non-degeneracy conditions:
(4.1)
$$_0๐ฆ^{}\text{ is trivial;}$$
(4.2)
$$P^{()}\text{ is irreducible}.$$
The first of these conditions can thought of as analogous to the condition on a multitype branching process that every initial condition leads to a non-zero probability of eventual extinction. The example given in the next paragraph suggests that it is not possible to dispense with some condition of this type. The second condition is probably not essential, but without it the proof given below would be considerably more complicated.
The following example shows how it is possible for the process to be non-endogenous even if $`2\rho =1`$. Let $`S=\{1,+1\}`$ and $`E=\{0,1\}`$ with $`\nu (0)=\nu (1)=1/2`$ and $`\mu (1)=\mu (+1)=1/2`$. Suppose that
(4.3)
$$\varphi (x_0,x_1,z)=\mathrm{๐}(z=0)x_0+\mathrm{๐}(z=1)x_1.$$
$`S^2S^{}=\{(1,+1),(+1,1)\}`$ and the transition matrix $`P^{()}`$ is $`1/2`$ times the identity matrix, so plainly $`2\rho =1`$. Furthermore $`\xi _{\mathrm{}}_0`$ is orthogonal to every subspace $`๐ฆ_n`$ and hence to $`๐ฆ`$, in particular this shows that endogeny does not hold. However this example does not satisfy the strong symmetry condition we have assumed for $`\varphi `$, namely that $`\varphi (x_0,x_1,z)=\varphi (x_1,x_0,z)`$ for all $`x_0,x_1S`$ and $`zE`$. I do not know whether there are any examples, with $`S`$ finite, and this symmetry assumption upheld, but for which (4.1) fails.
The sequence of subspaces $`_0๐ฆ_n^{}`$ is decreasing, and since $`_0`$ is finite-dimensional, (4.1) can only hold if there exits some $`m`$ for which
(4.4)
$$_0๐ฆ_m^{}\text{ is trivial.}$$
Now consider the quadratic form $`(f,P_{๐ฆ_m}f)`$. If the preceding condition holds then this quadratic form restricted to $`f_0`$ is positive definite, and using again the fact that $`_0`$ is finite-dimensional we deduce that there exists an $`ฯต>0`$ such that
(4.5)
$$(f,P_{๐ฆ_m}f)ฯต(f,f)\text{ for all }f_0.$$
This extends, see the final paragraph of this section, to
(4.6)
$$(f,P_{๐ฆ_{n+m}}f)ฯต^{|S|}(f,f)\text{ for all }f_S,$$
for a subset $`SG_n`$, where $`n1`$ is arbitrary, whilst $`m`$ is as above. Decomposing $`f_n`$ into its components in the subspaces $`_S`$ as $`S`$ varies through subsets of $`G_n`$, and using (4.14) below, gives,
(4.7)
$$(f,P_{๐ฆ_{n+m}}f)\underset{SG_n}{}ฯต^{|S|}(f,P__Sf)=(f,P_{๐ฆ_n}f)+\underset{\begin{array}{c}SG_n\\ S\mathrm{}\end{array}}{}ฯต^{|S|}(f,P__Sf)$$
Consequently for a fixed $`f`$,
(4.8)
$$\underset{\begin{array}{c}SG_n\\ S\mathrm{}\end{array}}{}ฯต^{|S|}(f,P__Sf)0,$$
as $`n\mathrm{}`$.
The criticality assumption that $`2\rho =1`$ implies, in the presence of the additional condition (4.2), that the sequence of matrices $`2^nP_n^{()}`$ is bounded as $`n`$ varies. Consequently, for a fixed $`f`$,
(4.9)
$$(f,A_nf)=\underset{SG_n}{}|S|(f,P__Sf)$$
is also bounded. The only way that this is consistent with (4.8) is for
(4.10)
$$\underset{\begin{array}{c}SG_n\\ S\mathrm{}\end{array}}{}(f,P__Sf)0,$$
which proves endogeny.
A couple of the steps used above need amplification. Start by considering a generalization of (2.1) and (2.2). By decomposing the tree at level $`n`$ we obtain a natural isomorphism
(4.11)
$$_{n+m}^{G_n}L^2(\mathrm{\Omega }_m,_m,m_m)^{\mathrm{\Gamma }_{n1}}L^2(E,,\nu ).$$
Splitting $`L^2(\mathrm{\Omega }_m,_m,m_m)`$ into the space of constants together its orthogonal complement $`L_0^2(\mathrm{\Omega }_m,_m,m_m)`$, we obtain the decomposition
(4.12)
$$_{n+m}=\underset{SG_n}{}_S(m),$$
where for each subset $`S`$ the subspace $`_S(m)`$ is generated by vectors of the form $`_{uG_n}f_u_{v\mathrm{\Gamma }_{n1}}g_v`$ with $`f_uL_0^2(\mathrm{\Omega }_m,_m,m_m)`$ for $`uS`$, $`f_u`$ being a constant for $`uS`$, and $`g_vL^2(E,,\nu )`$ being unrestricted. Notice that for each subset $`S`$ of $`G_n`$ the subspace $`_S`$ is included in $`_S(m)`$. There is a corresponding decomposition
(4.13)
$$๐ฆ_{n+m}=\underset{SG_n}{}๐ฆ_S(m),$$
where $`๐ฆ_S(m)`$ is a subspace of $`_S(m)`$. The orthogonal projection $`P_{๐ฆ_{n+m}}`$ acts on $`_{n+m}`$ by projecting each subspace $`_S(m)`$ onto the corresponding subspace $`๐ฆ_S(m)`$. Accordingly if $`f_{n+m}`$ is decomposed as $`f=f_S`$ with $`f_S_S(m)`$ then
(4.14)
$$(f,P_{๐ฆ_{n+m}}f)=\underset{S}{}(f_S,P_{๐ฆ_{n+m}}f_S)$$
a fact that was used at (4.7).
Further examination reveals that $`_S(m)`$ is naturally isomorphic to a tensor product of $`๐ฆ_n`$ together with $`|S|`$ copies of $`_m^0`$, where the latter is the orthogonal complement of the space of constants in $`_m`$. Similarly $`๐ฆ_S(m)`$ is naturally isomorphic to a tensor product of $`๐ฆ_n`$ together with $`|S|`$ copies of $`๐ฆ_m^0`$. The restriction of $`P_{๐ฆ_{n+m}}`$ to $`_S(m)`$ respects this tensor product structure, acting as the identity on the factor of $`๐ฆ_n`$ tensored with copies of the natural projection from $`_m^0`$ to $`๐ฆ_m^0`$ on the other factors. Similarly the restriction of $`P__nP_{๐ฆ_{n+m}}`$ to $`_S(m)`$ is the tensor product of the identity on $`๐ฆ_n`$ with $`|S|`$ copies of $`P__0P_{๐ฆ_m}`$ acting on $`_m^0`$. The inequality (4.6) follows from this and (4.5) since the smallest eigenvalue of a tensor product of operators is the product of the smallest eigenvalues.
## 5 Convergence to the dynamics
Throughout this section we work with the case $`2\rho >1`$ and we make the additional assumption
(5.1)
$$P^{()}\text{ is a primitive matrix}.$$
According to Perron-Frobenius theory, under this condition, the limit of the rescaled $`n`$-step transition matrices $`\rho ^nP_n^{()}`$ exists and is given by
(5.2)
$$\underset{n\mathrm{}}{lim}\rho ^nP_n^{()}((x,x^{}),(y,y^{}))=\theta (x,x^{})\kappa (y,y^{}),$$
where $`\theta `$ is the left eigenvector of $`P^{()}`$ corresponding to $`\rho `$, and, as before, $`\kappa `$ is the right eigenvector. Here we normalize $`\theta `$ and $`\kappa `$ so that
(5.3)
$$\underset{(x,x^{})S^2S^{}}{}\theta (x,x^{})\kappa (x,x^{})=1.$$
We may also normalize so that
(5.4)
$$\underset{(x,x^{})S^2S^{}}{}\theta (x,x^{})\mu (x)\mu (x^{})=1,$$
which fixes a choice of $`\kappa `$. We assume throughout this section that $`Q`$ is defined by (3.5) with this choice of $`\kappa `$. We deduce that
(5.5)
$$\rho ^n(\mu \mu \mu ^{})P_n^{(2)}\kappa ^{},$$
where $`\kappa ^{}(x,x^{})=\kappa (x,x^{})`$ when $`xx^{}`$ and $`\kappa ^{}(x,x)=_{x^{}x}\kappa (x,x^{})`$. Using the relationship (3.1) between the superoperator $`๐ซ`$ and $`P^{(2)}`$ this may be recast as
(5.6)
$$\underset{n\mathrm{}}{lim}\rho ^n๐ซ^nP_\mathrm{๐}^{}=Q.$$
Suppose that $`f,g_m`$ then a straightforward application of Proposition 4 gives
(5.7)
$$\frac{1}{(2\rho )^n}(f,A_ng)=\frac{1}{2^m\rho ^n}\underset{uG_m}{}(f,(๐ซ^{nm}P_\mathrm{๐}^{})^{(u)}g).$$
Thus in view of (5.6) we deduce that
(5.8)
$$\frac{1}{(2\rho )^n}(f,A_ng)(f,Q_mg)=(f,g),\text{ for all }f,g_m.$$
Recall that we proved that the form $``$ is closable, and let $`Q_{\mathrm{}}`$ be the self-adjoint operator associated with its closure. If $`f_m_m`$ the limit $`\stackrel{ฬ}{Q}_{\mathrm{}}f=lim_nQ_nf`$ exists and defines a operator $`\stackrel{ฬ}{Q}_{\mathrm{}}`$ with domain $`_m`$. The self -adjoint operator $`Q_{\mathrm{}}`$ is an extension ( the Friedrichs extension) of $`\stackrel{ฬ}{Q}_{\mathrm{}}.`$
###### Proposition 7.
If $`f_m`$ for some $`m`$, then as $`n`$ tends to infinity
$$(2\rho )^nA_nfQ_{\mathrm{}}f,$$
in the metric topology of $``$.
###### Proof.
Convergence in the metric topology is implied by weak convergence together with convergence of the norms. Consequently it is sufficient to verify that if $`f_m`$ for some $`m`$, then as $`n`$ tends to infinity,
$$(2\rho )^nA_nfQ_{\mathrm{}}f,$$
noting that weak convergence follows from this and (5.8).
We begin by computing $`Q_{\mathrm{}}f`$. We have, for $`f_m`$, and $`n>m`$,
$$\begin{array}{c}(f,Q_n^2f)=(2\rho )^{2n}\underset{uG_n}{}(f,(Q^2)^{(u)}f)+(2\rho )^{2n}\underset{\begin{array}{c}u,vG_n\\ uv\end{array}}{}(f,Q^{(u)}Q^{(v)}f)\hfill \\ \hfill =(2\rho )^{2n}\underset{uG_n}{}(f,(Q^2)^{(u)}f)+2\underset{r=m}{\overset{n1}{}}2^{2(nr1)}(2\rho )^{2n}\underset{uG_r}{}(f,(๐ฌ๐ซ^{nr1}Q)^{(u)}f)\\ \hfill +2^{2(nm)}(2\rho )^{2n}\underset{\begin{array}{c}u,vG_m\\ uv\end{array}}{}(f,(๐ซ^{nm}Q)^{(u)}(๐ซ^{nm}Q)^{(v)}f)\\ \hfill =(2\rho )^{2n}\underset{uG_n}{}(f,(Q^2)^{(u)}f)+2\underset{r=m}{\overset{n1}{}}(2\rho )^{2(r+1)}\underset{uG_r}{}(f,(๐ฌQ)^{(u)}f)\\ \hfill +(2\rho )^{2m}\underset{\begin{array}{c}u,vG_m\\ uv\end{array}}{}(f,Q^{(u)}Q^{(v)}f).\end{array}$$
Now as $`n`$ tends to infinity $`Q_nf`$ converges to $`Q_{\mathrm{}}f`$ in the metric topology of $``$, and hence the limit of the lefthandside above is $`Q_{\mathrm{}}f^2`$. Turning to the righthandside the first term tends to zero, and consequently we deduce that
$$Q_{\mathrm{}}f^2=(2\rho )^{2m}\underset{\begin{array}{c}u,vG_m\\ uv\end{array}}{}(f,Q^{(u)}Q^{(v)}f)+2\underset{r=m}{\overset{\mathrm{}}{}}(2\rho )^{2r2}\underset{uG_r}{}(f,(๐ฌQ)^{(u)}f).$$
A similar calculation is valid for $`A_nf^2`$. If we denote $`P_\mathrm{๐}^{}`$ by $`L`$ and $`\rho ^r๐ซ^rP_\mathrm{๐}^{}`$ by $`L_r`$, then
$$\begin{array}{c}(2\rho )^{2n}A_nf^2=(2\rho )^{2n}\underset{uG_n}{}(f,L^{(u)}f)+2\underset{r=m}{\overset{n1}{}}(2\rho )^{2r2}\underset{uG_r}{}(f,(๐ฌL_{nr1})^{(u)}f)\hfill \\ \hfill +(2\rho )^{2m}\underset{\begin{array}{c}u,vG_m\\ uv\end{array}}{}(f,L_{nm}^{(u)}L_{nm}^{(v)}f).\end{array}$$
Letting $`n`$ tend to infinity we observe that since $`L_n`$ tends to $`Q`$ we obtain termwise convergence to the expression for $`Q_{\mathrm{}}f^2`$. To complete the proof we appeal to dominated convergence noting that since $`๐ฌL_n1=0,`$ we have the estimate
$$\underset{uG_r}{}(f,(๐ฌL_{nr1})^{(u)}f)K(f,A_rf),$$
where $`K=sup_n๐ฌL_n<\mathrm{}`$. โ
###### Lemma 8.
For $`f_n`$,
$$Q_{\mathrm{}}f\frac{K}{(2\rho )^n}A_nf,$$
where $`K=sup\{Q_{\mathrm{}}f:f_0\text{ and }f=1\}`$.
###### Proof.
Recall the natural isomorphism
$$_n^{G_n}L^2(S,๐ฎ,\mu )^{\mathrm{\Gamma }_{n1}}L^2(E,,\nu ).$$
Similarly
$$^{G_n}L^2(\mathrm{\Omega },,m)^{\mathrm{\Gamma }_{n1}}L^2(E,,\nu ).$$
The operator $`Q_{\mathrm{}}:_n`$ can be written as a sum
$$Q_{\mathrm{}}=\frac{1}{(2\rho )^n}\underset{uG_n}{}Q_{\mathrm{}}^{(u)},$$
where $`Q_{\mathrm{}}^{(u)}`$ is unitary equivalent the tensor product of $`Q_{\mathrm{}}:L^2(S,๐ฎ,\mu )L^2(\mathrm{\Omega },,m)`$ on the factor corresponding to the node $`uG_n`$ and the identity on all other factors. It is easy to see from this structure, together with the fact that $`Q_{\mathrm{}}1=0`$ that if $`f_n`$ decomposes as $`f=_{SG_n}f_S`$ with $`f_S_S`$ then
$$\begin{array}{c}Q_{\mathrm{}}f^2=\underset{S}{}Q_{\mathrm{}}f_S^2=\frac{1}{(2\rho )^{2n}}\underset{S}{}\underset{u,vS}{}(Q_{\mathrm{}}^{(u)}f_S,Q_{\mathrm{}}^{(v)}f_S)\hfill \\ \hfill \frac{K^2}{(2\rho )^{2n}}\underset{S}{}|S|^2f_S^2=\frac{K^2}{(2\rho )^{2n}}A_nf^2.\end{array}$$
###### Lemma 9.
$`_m_m`$ is a core for $`Q_{\mathrm{}}`$.
###### Proof.
Recall that $`\stackrel{ฬ}{Q}_{\mathrm{}}`$ is the restriction of $`Q_{\mathrm{}}`$ to $`_m`$. To verify the claimed result it suffices to show that, for some $`\alpha >0`$, the range of $`(\alpha \stackrel{ฬ}{Q}_{\mathrm{}})`$ is dense in $``$. For this shows that $`\alpha `$ belongs to the resolvent set of the closure of $`\stackrel{ฬ}{Q}_{\mathrm{}}`$, and then we apply the criterion of Theorem X.1 of .
Fix some $`\alpha >0`$. Let $`R_n^\alpha :_n_n`$ be the $`\alpha `$ resolvent of $`Q_n`$. Given $`f_m_m`$ let $`v_n=R_n^\alpha f`$ for $`n`$ sufficiently large. Let $`g_m`$, then for all sufficiently large $`n`$,
$$(g,(\alpha Q_{\mathrm{}})v_n)=(g,(\alpha Q_n)v_n)=(g,f).$$
Now suppose that we know that $`(\alpha Q_{\mathrm{}})v_n`$ is uniformly bounded in norm, then we deduce that, for any $`g`$,
$$(g,(\alpha Q_{\mathrm{}})v_n)(g,f)\text{ as }n\mathrm{}.$$
Thus the range of $`(\alpha \stackrel{ฬ}{Q}_{\mathrm{}})`$ is weakly dense, and consequently norm dense in $``$.
To verify the supposition that $`(\alpha Q_{\mathrm{}})v_n`$ is uniformly bounded in norm we note that since $`v_n`$ are uniformly bounded in suffices to verify that $`Q_{\mathrm{}}v_n`$ are also. Using the previous lemma and the fact that $`R_n^\alpha `$ and $`A_n`$ commute we have
$$\begin{array}{c}Q_{\mathrm{}}v_n\frac{K}{(2\rho )^n}A_nv_n=\frac{K}{(2\rho )^n}A_nR_n^\alpha f=\frac{K}{(2\rho )^n}R_n^\alpha A_nf\frac{K}{\alpha (2\rho )^n}A_nf,\hfill \end{array}$$
and the righthandside is bounded as $`n`$ tends to infinity. โ
###### Proof of Theorem 2.
We know that $`_n`$ is a common core for $`Q_{\mathrm{}}`$ and for $`A_nP__n`$. Thus the convergence established at Proposition 7 implies strong resolvent convergence of $`(2\rho )^nA_nP__n`$ to $`Q_{\mathrm{}}`$ (see Theorem VIII.25 of ). โ
Convergence of the generators in the strong resolvent sense implies that the semigroups converge in the strong operator topology. From this fact we obtain the following corollary, which can can also be expressed in terms of $`T^{(2)}`$.
###### Corollary 10.
For each fixed $`t>0`$, as $`n\mathrm{}`$,
$$๐ฌ^n\left(P_\mathrm{๐}+e^{(2\rho )^nt}P_\mathrm{๐}^{}\right)M_t,$$
where $`M_t`$ acting on $`L^2(S)`$ is defined by identifying $`_0`$ with $`L^2(S)`$ and then setting $`(f,M_tg)_{L^2(S)}=(f,e^{tQ_{\mathrm{}}}g)`$.
Another interpretation of this corollary is available in terms of the spectral measures $`\mu _f^{(n)}`$ defined at (1.10). The rescaled measures
(5.9)
$$\stackrel{~}{\mu }_f^{(n)}\left([0,x]\right)=\mu _f^{(n)}(([0,(2\rho )^nx])\text{ for }x0,$$
converge weakly towards a measure $`\mu _f`$ whose Laplace transform is given by
(5.10)
$$_0^{\mathrm{}}e^{tx}\mu _f(dx)=(f,e^{tQ_{\mathrm{}}}f).$$
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# On the Throughput-Delay Tradeoff in Cellular Multicast
## 1 Introduction
Traditional information theoretic investigations pay little, if any, attention to the notion of delay. Clearly, this approach is not adequate for many applications, especially those with strict Quality of Service (QoS) constraints. To avoid this shortcoming, recent years have witnessed a growing interest in cross layer design approaches. The underlying idea in these approaches is to jointly optimize the physical, data link, and networking layers in order to satisfy the QoS constraints with the minimum expenditure of network resources. Early investigations on cross layer design have focused on the single user case . These works have shed light on the fundamental tradeoffs in this scenario and devised efficient power and rate control policies that approach these limits. More recent works have considered multi-user cellular networks . These studies have enhanced our understanding of the fundamental limits and the structure of optimal resource allocation strategies. Here, we take a first step towards generalizing this cross layer approach to the wireless multicast scenario. This scenario is characterized by a strong interaction between the network, medium access, and physical layers. This interaction adds significant complexity to the problem which motivated the adoption of a simplified on-off model for the wireless channel in several of the recent works on wireless multicast . In the sequel, we argue that employing more accurate models for the wireless channel allows for valuable opportunities for exploiting the wireless medium to yield performance gains. More specifically, our work sheds light on the role of the following characteristics of the wireless channel in the design of multicast scheduling strategies: 1) The multi-user diversity resulting from the statistically independent channels seen by the different users , 2) The wireless multicast gain resulting from the fact that any information transmitted over the wireless channel is overheard by all users, possibly with different attenuation factors, and 3) The cooperative gain resulting from antenna sharing between users .
To illustrate the main ideas, we first focus on the single group (pure multicast) scenario where the same information stream is transmitted to all users in the network . We consider three classes of scheduling algorithms with progressively increasing complexity. The first class strives for minimum complexity by resorting to a static scheduling strategy along with memoryless decoding<sup>1</sup><sup>1</sup>1Memoryless decoding refers to the fact that the decoder memory is flushed in case of decoding failure.. In this approach, we schedule transmission to a fraction of the users that enjoy favorable channel conditions. While the identity of the target users change, based on the channel conditions, the static nature of the algorithm is manifested in the fact that a fixed fraction of the users is able to decode every transmitted packet. We establish the throughput-delay tradeoff allowed by varying the fraction of users targeted in every transmission. To gain more insight into the problem, we study in more detail the three special cases of scheduling transmissions to the best, worst and median user<sup>2</sup><sup>2</sup>2These notions will be defined rigorously in the sequel.. Here we establish the asymptotic throughput optimality of the median user scheduler and show that the price for this optimality is an exponential growth in delay with the number of users. The second scheduling policy resorts to a higher complexity incremental redundancy encoding/decoding strategy to achieve a better throughput-delay tradeoff. This scheme is based on a hybrid Automatic Repeat reQuest (ARQ) strategy and is shown to yield a significant reduction in the delay, compared with the median user scheduler, at the expense of a minimal penalty in the throughput. The third, and most complex, scheduling strategy benefits from the cooperation between the different users to minimize the delay while achieving the optimal scaling law of the throughput. More specifically, we show that the proposed cooperative multicast strategy simultaneously achieves the optimal scaling laws of both throughput and delay at the expense of a high complexity. Finally, we extend our study to the multi-group scenario where independent streams of information are transmitted to different groups of users. Here, we generalize our scheduling algorithms to exploit the multi-group diversity available in such scenarios.
The rest of the paper is organized as follows. In Section 2, we introduce the system model along with our notation. In Section 3, we propose the three classes of scheduling algorithms for the pure multicast scenario and characterize the achieved throughput-delay tradeoffs. We then extend our schemes to exploit the multi-group diversity in Section 4. The potential performance gains allowed by multi-transmit antenna base stations are quantified in Section 5. In Section 6, we present numerical results that validate our theoretical claims in certain representative scenarios. Finally, some concluding remarks are offered in Section 7. In order to enhance the flow of the paper, we collect all the proofs in the Appendices.
## 2 System Model
We consider the downlink of a single cell system where a base station serves $`G`$ groups of users. The information streams requested by the different groups from the base station are independent of each other. Each group consists of $`N`$ users. All the users within a group request the same information from the base station. Unless otherwise stated, the base station is assumed to be equipped with a single transmit antenna. Each user is assumed to have only a single receive antenna. We consider time-slotted transmission in which the signal received by user $`i`$ in time slot $`k`$ is given by
$$y_i[k]=h_ix[k]+n_i[k],$$
where $`x[k]`$ denotes the complex-valued signal transmitted by the base station in slot $`k`$, $`h_i`$ represents the complex flat fading coefficient of the channel between the base station and the $`i^{th}`$ user, and $`n_i[k]`$ represents the zero-mean unit-variance complex additive white Gaussian noise at the $`i^{th}`$ user in slot $`k`$. The noise processes are assumed to be circularly symmetric and independent across users. The channel between the base station and each user is assumed to be quasi-static with coherence time $`T_c`$. Thus the fading coefficients remain constant throughout an interval of length $`T_c`$ and change independently from one interval to the next. The fading coefficients $`\{h_i\}`$ are assumed to be independent and identically distributed (i.i.d.) across the users and follow a Rayleigh distribution with $`E\left[|h_i|^2\right]=1`$. In this paper, we restrict our attention to this symmetric scenario, and hence, issues related to fairness are outside the scope of this work. Each packet transmitted by the base station is assumed to be of constant size $`S`$. We further employ the following short term average power constraint
$$E\left[|x[k]|^2\right]P.$$
Clearly, further performance gain may be reaped through a carefully constructed power allocation policy if this short term power constraint is replaced by a long term one. This line of work, however, is not pursued here and we only rely on rate adaptation and scheduling based on the instantaneous channel state. The scheduling schemes proposed in the sequel require one further assumption. We require all the channel gains to be available at the base station. Hence the proposed scheduling strategies, except the incremental redundancy scheme<sup>3</sup><sup>3</sup>3For the incremental redundancy scheme, the base station only needs to know when to stop transmission of the current codeword., assume perfect knowledge of the channel state information (CSI) at both the transmitter and receiver. In our throughput analysis, we use capacity expressions for the channel transmission rates. Here we implicitly assume that the base station employs coding schemes that approach the channel capacity which justify our use of the fundamental information theoretic limit of the channel.
In our delay analysis, we consider backlogged queues, and hence, the only meaningful measure of delay is the transmission delay. This leads to the following definitions for throughput and delay that will be adopted in the sequel.
###### Definition 1
The throughput of a scheduling scheme is defined as the sum of the throughputs provided by the base station to all the individual users within all the groups in the system.
###### Definition 2
The delay of a scheduling scheme is defined as the delay between the instant representing the start of transmission of a packet belonging to a particular group of users, and the instant when the packet is successfully decoded by all the users in that group.
A brief comment on the notion of delay adopted in our work is now in order. This definition suffers from the fundamental weakness that it does not account for the queuing delay experienced by the packets. Unfortunately, at the moment we do not have an analytical characterization of the queuing delay for the general case. However, as argued in the sequel, our delay analysis offers a lower bound on the total delay which is very tight in several important special cases. Furthermore, this analysis provides a very useful tool for rank-ordering the different classes of scheduling algorithms and sheds light on their structural properties.
To facilitate analytical tractability, we focus on evaluating the asymptotic scaling laws of the throughput and delay in the sequel. In this analysis, we use the following set of Knuthโs asymptotic notations throughout the paper: 1)$`f(n)=O(g(n))`$ iff there are constants $`c`$ and $`n_0`$ such that $`f(n)cg(n)`$ $`n>n_0`$, 2) $`f(n)=\mathrm{\Omega }(g(n))`$ iff there are constants $`c`$ and $`n_0`$ such that $`f(n)cg(n)`$ $`n>n_0`$, and 3) $`f(n)=\mathrm{\Theta }(g(n))`$ iff there are constants $`c_1`$, $`c_2`$ and $`n_0`$ such that $`c_1g(n)f(n)c_2g(n)`$ $`n>n_0`$. Furthermore, the two following technical assumptions are imposed.
1. We let
$$T_c=\mathrm{\Theta }\left(\frac{1}{\mathrm{log}\mathrm{log}NG}\right).$$
(1)
This technical assumption is made to ensure (as shown in the sequel) that the average service time required for transmitting a packet is not dominated by the scaling behavior of $`T_c`$.
2. In our delay analysis, we make an exponential server assumption, i.e., the rate of service $`R`$ offered by the server in any time slot is assumed to follow an exponential distribution with the same mean as that obtained from our problem formulation. Thus, for a particular scheduling algorithm, the service rate distribution is given by
$$F_R(r)=1e^{\mu r},r0,$$
(2)
where $`\mu =(1/E[R])`$ depends on the channel characteristics and the scheduling algorithm.
## 3 Single Group (Pure Multicast) Scenario
In this section, we consider the pure multicast scenario where the same information stream is transmitted to all users in the network. In the non-cooperative scenario, the throughput-optimal scheme is an $`N`$-level superposition coding/successive decoding scheme . This strategy, however, suffers from excessive complexity and the corresponding delay analysis seems intractable at the moment. This motivates our work where we focus on the throughput-delay tradeoff of low complexity scheduling schemes. Interestingly, we identify a low complexity static scheduling scheme, as defined in the next section, that achieves the optimal scaling law of the throughput. Furthermore, we establish the optimality of the proposed cooperative multicast scheme in terms of the scaling laws of both delay and throughput.
### 3.1 Static Scheduling With Memoryless Decoding
In this class of scheduling algorithms, referred to as static schedulers in the sequel, we schedule transmission to a fixed fraction of the users with favorable channel conditions. The transmission rate is adjusted such that each transmission by the base station is intended for successful reception by $`(N/\alpha )`$ users in the system. Hence at any time instant, the base station transmits to the user whose instantaneous SNR occupies the $`(N(N/\alpha )+1)^{th}`$ position in the ordered list of instantaneous SNRs of all users. The other $`((N/\alpha )1)`$ users with higher channel gains can also decode the transmitted information. The parameter $`\alpha `$ of the scheme is restricted to be a factor of $`N`$ and satisfies $`\alpha ๐ต^+`$ and $`1\alpha N`$. This scheme is โstaticโ in the sense that the fraction of users targeted in every transmission remains the same (i.e., the parameter $`\alpha `$ is not a function of time). When $`\alpha >1`$, some of the users will not be able to decode. The memoryless property dictates that those users flush their memories and wait for future re-transmissions of the packet. This assumption is imposed to limit the complexity of the encoding/decoding process. In Section 3.2, we relax this memoryless decoding assumption and quantify the gains offered by carefully constructed ARQ schemes. As shown later, this class of static scheduling algorithms exploit both the multi-user diversity and multicast gains, to varying degrees, depending on the parameter $`\alpha `$.
The average throughput of this general static scheduling scheme is given by
$$R_{tot}=\left(\frac{N}{\alpha }\right)E[R_\alpha ],$$
where $`R_\alpha `$ is the transmission rate to each of the intended $`(N/\alpha )`$ users and is given by
$$R_\alpha =\mathrm{log}\left(1+|h_{\pi (N\frac{N}{\alpha }+1)}|^2P\right),$$
(3)
where $`|h_{\pi (N\frac{N}{\alpha }+1)}|^2`$ is the channel power gain of the user whose SNR occupies the $`(N(N/\alpha )+1)^{th}`$ position in the ordered list of SNRs of all users. Throughout the paper, the $`\mathrm{log}(.)`$ function refers to the natural logarithm, and hence, the average throughput is expressed in nats.
A critical step in the delay analysis is to identify the queuing model. In our model, the base station maintains $`\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)`$ queues, one for each combination of $`(N/\alpha )`$ users. These queues can be divided into sets with $`\alpha `$ coupled queues in each set such that the combinations of users served by the $`\alpha `$ queues within a set are mutually exclusive (to ensure that multiple copies of the same packet are not sent to any of the users) and collectively exhaustive (to ensure that the packet reaches all the users), i.e., every user in the system is served by exactly one of the $`\alpha `$ queues in each set. For example, with $`N=6`$ users and $`\alpha =3`$, we have $`15`$ queues divided into $`5`$ sets with three queues in each set (One possible set of coupled queues serve users $`\{(1,2),(3,4),(5,6)\}`$ and another possible set may serve users $`\{(1,4),(2,5),(3,6)\}`$. Note that each user occurs once and only once in each set). Hence, any packet that arrives at the base station is routed towards one of the sets<sup>4</sup><sup>4</sup>4Here, we use a probabilistic approach for choosing the set with a uniform distribution. where it is stored in all the $`\alpha `$ queues within that set (since it needs to be transmitted to all the users in the system). Thus the delay in transmitting a particular packet to all the users is given by the delay in transmitting that packet from each of the $`\alpha `$ coupled queues in the corresponding set. Moreover, the base station services only one of the $`\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)`$ queues at any time, which is chosen based on the instantaneous fading coefficients of all the users. An example of the queuing model for a system with $`N=6`$ users and $`\alpha =3`$ is shown in Fig. 1.
In our analysis, we benefit from the concept of worst case delay proposed in for analyzing the delay in unicast networks. In this work, the authors characterized the worst case delay by restating their problem as the โcoupon collector problemโ which has been studied extensively in the mathematics literature . In the coupon collector problem, the users are assumed to have coupons and the transmitter is the collector that selects one of the users randomly (with uniform distribution) and collects his coupon. The problem is to characterize the average number of trials required to ensure that the collector collects $`m`$ coupons from all the users. Our queuing problem is analogous to the coupon collector problem with the only fundamental difference being that the size of the coupons is time-varying in our problem due to rate adaptation (the detailed analysis is presented in the proofs). Now, we are ready to state our result that characterizes the scaling laws of throughput and delay for the different static scheduling algorithms.
###### Theorem 3
The average throughput $`R_{tot}`$ of the general static scheduling scheme is given by
$$R_{tot}=\frac{N}{\alpha }\left[\underset{i=1}{\overset{N}{}}\left(\genfrac{}{}{0pt}{}{N}{i}\right)(1)^ie^{\left(\frac{i}{P}\right)}Ei\left(\frac{i}{P}\right)+\underset{k=(N\frac{N}{\alpha }+1)}{\overset{N1}{}}\left(\genfrac{}{}{0pt}{}{N}{k}\right)\left[\underset{i=0}{\overset{k}{}}\left(\genfrac{}{}{0pt}{}{k}{i}\right)(1)^ie^{\left(\frac{Nk+i}{P}\right)}Ei\left(\frac{(Nk+i)}{P}\right)\right]\right],$$
(4)
where
$$Ei(x)=_{\mathrm{}}^x\frac{e^t}{t}dt.$$
The average delay $`D`$ of this scheme satisfies
$$D=\mathrm{max}\{\mathrm{\Omega }\left(\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)\frac{\mathrm{log}\alpha }{\mathrm{log}\mathrm{log}N}\right),\mathrm{\Omega }\left(\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)E[X_{min}]\right)\},$$
(5)
where $`X_{min}=\mathrm{min}_{i=1}^\alpha X_i`$ and the $`X_i`$โs are defined as the service times required for transmitting a packet from the $`i^{th}`$ queue of a set of $`\alpha `$ queues assuming that the server always services the $`i^{th}`$ queue.
To gain more insights into the rather involved throughput and delay expressions of Theorem 3, we study three special cases of the general static scheduling scheme in more detail. This detailed analysis sheds light on the throughput-delay tradeoff achievable by varying $`\alpha `$. We further establish the optimality of the scheduler corresponding to $`\alpha =2`$ with respect to the throughput scaling law.
#### 3.1.1 Worst User Scheduler
The worst user scheme corresponds to the case $`\alpha =1`$ of the general scheduling scheme. This scheme maximally exploits the multicast gain by always transmitting to the user with the least instantaneous SNR. This enables all the users to successfully decode the transmission and thus any particular packet reaches all the users in a single transmission. However, the multi-user diversity inherent in the system works against the performance of this scheme and results in a decrease in the individual throughput to any user.
The average throughput of the worst user scheme is given by
$$R_{tot}=NE\left[\mathrm{log}\left(1+|h_{\pi (1)}|^2P\right)\right],$$
where $`|h_{\pi (1)}|^2`$ is the minimum channel gain among all the $`N`$ users in the system, whose distribution and density functions are given by
$$F_{|h_{\pi (1)}|^2}(x)=1e^{Nx}\text{and}f_{|h_{\pi (1)}|^2}(x)=Ne^{Nx},x0.$$
For implementing this scheme, the base station needs to maintain only a single queue that caters to all the users in the system.
###### Lemma 4
The average throughput of the worst user scheme scales as
$$R_{tot}=\mathrm{\Theta }(1)$$
(6)
with the number of users $`N`$. The average delay of this scheme scales as
$$D=\mathrm{\Theta }(N).$$
(7)
Thus the average throughput of the worst user scheme does not scale with the number of users $`N`$ while the average delay increases linearly with $`N`$.
#### 3.1.2 Best User Scheduler
This scheme corresponds to the case $`\alpha =N`$ of the general scheduling scheme and maximally exploits the multi-user diversity available in the system. Since the transmission rate is adjusted based on the user with the maximum instantaneous SNR, this scheme fails to exploit any of the multicast gain and any particular packet must be repeated $`N`$ times. The average throughput of the best user scheme is given by
$$R_{tot}=E\left[\mathrm{log}\left(1+|h_{\pi (N)}|^2P\right)\right],$$
where $`|h_{\pi (N)}|^2`$ is the maximum channel gain among all the $`N`$ users in the system, whose distribution function is given by
$$F_{|h_{\pi (N)}|^2}(x)=\left(1e^x\right)^N,x0.$$
In this special case, the base station maintains $`N`$ queues, one for each user in the system, and any packet that arrives into the system enters all the $`N`$ queues. The following result establishes the throughput and delay scaling laws achieved by the best user scheduler.
###### Lemma 5
The average throughput of the best user scheme scales as
$$R_{tot}=\mathrm{\Theta }(\mathrm{log}\mathrm{log}N)$$
(8)
with the number of users $`N`$. The average delay of this scheme scales as
$$D=\mathrm{\Omega }\left(\frac{N\mathrm{log}N}{\mathrm{log}\mathrm{log}N}\right).$$
(9)
From Lemmas 4 and 5, one can conclude that maximally exploiting the multi-user diversity yields higher throughput gains than maximally exploiting the multicast gain. This throughput gain, however, is obtained at the expense of a higher delay. This observation motivates the investigation of other variants of the static scheduling strategy which achieve other points on the throughput-delay tradeoff.
#### 3.1.3 Median User Scheduler
The median user scheduler corresponds to the case $`\alpha =2`$ of the general scheduling scheme. In this scheme, the base station maintains $`\left(\genfrac{}{}{0pt}{}{N}{N/2}\right)`$ queues, one for each combination of $`(N/2)`$ users. This scheme strikes a balance between exploiting multi-user diversity and multicast gain. The base station always transmits to the user whose instantaneous SNR occupies the median position of the ordered list of SNRs. Each transmission is, therefore, successfully decoded by half the users in the system and the same information needs to be repeated only twice before it reaches all the users. Thus, unlike the best user scheduler, this scheduler benefits from the wireless multicast gain. Moreover, unlike the worst user scheduler, the inherent multi-user diversity does not degrade the performance of this scheduler (since the instantaneous SNR of the median user is not expected to degrade with $`N`$). In fact, we show in the following that this scheme achieves the optimal scaling law of the throughput as the number of users $`N`$ grows to infinity.
###### Lemma 6
The proposed median user scheme achieves the optimal scaling law of the throughput. The average throughput of this scheme scales as
$$R_{tot}=\mathrm{\Theta }(N)$$
(10)
with the number of users $`N`$. The average delay of this scheme scales as
$$D=\mathrm{\Theta }\left(\left(\genfrac{}{}{0pt}{}{N}{N/2}\right)\right)=\mathrm{\Theta }\left(\frac{2^N}{\sqrt{N}}\right).$$
(11)
Thus the throughput optimality of the median user scheduler is obtained at the expense of an exponentially increasing delay with the number of users $`N`$. Overall, these three special cases of the static scheduling strategy show that one can achieve different points on the throughput-delay tradeoff by varying $`\alpha `$.
### 3.2 Incremental Redundancy Multicast
In this section, we relax the memoryless decoding requirement and propose a scheme that employs a higher complexity incremental redundancy encoding/decoding strategy to achieve a better throughput-delay tradeoff than the static scheduling schemes. The proposed scheme is an extension of the incremental redundancy scheme given by Caire et al in . An information sequence of $`b`$ bits is encoded into a codeword of length $`LM`$, where $`M`$ refers to the rate constraint. The first $`L`$ bits of the codeword are transmitted in the first attempt. If a user is unable to successfully decode the transmission, it sends back an ARQ request to the base station. If the base station receives an ARQ request from any of the users, it transmits the next $`L`$ bits of the same codeword in the next attempt. This process continues until either all $`N`$ users successfully decode the information or the rate constraint $`M`$ is violated. Then the codeword corresponding to the next $`b`$ information bits is transmitted in the same fashion. In this scheme, even if some of the users successfully decode the information in very few attempts, they still have to wait until all the $`N`$ users successfully receive the information before any new information is transmitted to them by the base station. This sub-optimality of the proposed scheme results in significant complexity reduction by avoiding the use of superposition coding and successive decoding. Moreover, this scheme does not require the knowledge of perfect CSI at the base station. The base station only needs to know when to stop transmission of the current codeword. Hence the feedback required is minimal. The following result establishes the superior throughput-delay tradeoff achieved by this scheme, compared with the class of static schedulers with memoryless decoding.
###### Theorem 7
The average throughput of the incremental redundancy scheme scales as
$$R_{tot}=\mathrm{\Theta }\left(\frac{N\mathrm{log}\mathrm{log}N}{\mathrm{log}N}\right)$$
(12)
with the number of users $`N`$. The average delay $`D`$ of this scheme scales as
$$D=\mathrm{\Theta }\left(\frac{\mathrm{log}N}{\mathrm{log}\mathrm{log}N}\right).$$
(13)
Thus, we can see that incremental redundancy multicast avoids the exponentially growing delay of the median user scheduler at the expense of a minimal penalty in throughput. In fact, the loss in both delay and throughput scaling laws, compared to the optimal values, is only a factor of $`\mathrm{log}(N)/\mathrm{log}\mathrm{log}(N)`$. In this approach, the base station needs to maintain only a single queue that serves all the users in the system. This approach, however, entails added complexity in the incremental redundancy encoding and the storage and joint decoding of all the observations.
### 3.3 Cooperative Multicast
In this section, we demonstrate the benefits of user cooperation and quantify the tremendous gains that can be achieved by allowing the users to cooperate with each other. In particular, we propose a cooperation scheme that minimizes the delay while achieving the optimal scaling law of the throughput. This scheme is divided into two stages. In the first half of each time slot, the base station transmits the packet to one half of the users in the system (i.e., the median user scheduler). During the next half of the slot, the base station remains silent. Meanwhile all the users that successfully decoded the packet in the first half of the slot cooperate with each other and transmit the packet to the other $`(N/2)`$ users in the system. This is equivalent to a transmission from a transmitter equipped with $`(N/2)`$ transmit antennas to the worst user in a group of $`(N/2)`$ users. If $`R_{s1}`$ and $`R_{s2}`$ are the rates supported in the first and second stage respectively, then the actual transmission rate is chosen to be $`\mathrm{min}\{R_{s1},R_{s2}\}`$ in both stages of the cooperation scheme. Note that the rate $`R_{s2}`$ is chosen such that the information can be successfully decoded even by the worst of the remaining $`(N/2)`$ users. Here, we note that this scheme requires the base station to know the CSI of the inter-user channels. The scheme, however, does not require the users to have such transmitter CSI (i.e., in the second stage the users cooperate blindly by using i.i.d. random coding). The average throughput of the proposed cooperation scheme is thus given by
$$R_{tot}=\left(\frac{N}{2}\right)E\left[\mathrm{min}\{R_{s1},R_{s2}\}\right].$$
The following result establishes the optimality of the proposed scheme, in terms of both delay and throughput scaling laws.
###### Theorem 8
The proposed cooperation scheme achieves the optimal scaling laws of both delay and throughput. In particular, the average throughput of this scheme scales as
$$R_{tot}=\mathrm{\Theta }(N)$$
(14)
with the number of users $`N`$, while the average delay scales as
$$D=\mathrm{\Theta }(1).$$
(15)
Here we assume that the inter-user channels have the same fading statistics as the channels between the base station and users, and the total transmitted power is upper bounded by $`P`$.
The price for this optimal performance is the added complexity needed to 1) equip every user terminal with a transmitter, 2) decode/re-encode the information at each cooperating user terminal, and 3) inform the base station with perfect CSI of the inter-user channels.
## 4 Multi-group Diversity
In this section, we generalize the scheduling schemes proposed in Section 3 to the multi-group scenario where different information streams are requested by different subsets of the user population. We modify the proposed schemes to exploit the multi-group diversity available in this scenario by always transmitting to the best group. We characterize the asymptotic scaling laws of the throughput and delay of the static schedulers with the number of users per group $`N`$ and the number of groups $`G`$ in the following theorem.
###### Theorem 9
1. The average throughput of the best among worst users scheme scales as
$$R_{tot}=\mathrm{\Theta }(\mathrm{log}G)$$
(16)
with $`N`$ and $`G`$. The average delay of this scheme scales as
$$D=\mathrm{\Theta }\left(\frac{NG}{\mathrm{log}G}\right).$$
(17)
2. The average throughput of the best among best users scheme scales as
$$R_{tot}=\mathrm{\Theta }(\mathrm{log}\mathrm{log}NG)$$
(18)
with $`N`$ and $`G`$. The average delay of this scheme scales as
$$D=\mathrm{\Omega }\left(\frac{NG\mathrm{log}N}{\mathrm{log}\mathrm{log}NG}\right).$$
(19)
3. The average throughput of the best among median users scheme satisfies
$$\mathrm{\Omega }(N)=R_{tot}=O(N\mathrm{log}\mathrm{log}G),$$
(20)
while the average delay of this scheme satisfies
$$\mathrm{\Omega }\left(\frac{G2^N}{\sqrt{N}\mathrm{log}\mathrm{log}G}\right)=D=O\left(\frac{G2^N}{\sqrt{N}}\right).$$
(21)
In the multi-group incremental redundancy scheme, the information bits corresponding to each of the groups are encoded independently. During each time slot, the base station selects that group for which it can send the highest total instantaneous rate to the users who failed to decode up to this point. This selection process makes the scheme โdynamicโ in the sense that the outcome of the scheduling process at any particular time slot depends on the outcomes in all previous slots. Unfortunately, this dynamic nature of the proposed scheme adds significant complexity to the problem and, at the moment, we do not have an analytical characterization of the corresponding scaling laws.
In the multi-group cooperation scheme, during each time slot, the base station selects the best group $`\widehat{g}`$ for transmission according to the condition
$$\widehat{g}=\mathrm{arg}\underset{g=1,\mathrm{},G}{\mathrm{max}}\left\{\left(\frac{N}{2}\right)\mathrm{min}\{R_{s1}^g,R_{s2}^g\}\right\}.$$
(22)
###### Theorem 10
The average throughput of the proposed multi-group cooperation scheme satisfies
$$\mathrm{\Omega }(N)=R_{tot}=O\left(N\mathrm{log}\mathrm{log}G\right),$$
(23)
while the average delay of this scheme satisfies
$$\mathrm{\Omega }\left(\frac{G}{\mathrm{log}\mathrm{log}G}\right)=D=O(G).$$
(24)
As expected, the throughput gain resulting from the multi-group diversity entails a corresponding price in the increased delay.
## 5 Multi-Transmit Antenna Gain
The performance of the proposed static scheduling schemes depends on the spread of the fading distribution. For exploiting significant multi-user diversity gains, the distribution needs to be well-spread out. The lower the spread of the distribution, the lesser the multi-user diversity gain (or loss as shown in the following). To illustrate this point, we consider a scenario where the base station is equipped with $`L`$ transmit antennas. We assume that the base station has knowledge of only the total effective SNR at any particular user and does not know the individual channel gains from each transmit antenna to that user. Under this assumption, the base station just distributes the available power equally among all the $`L`$ transmit antennas. Thus the effective fading power gains follow a normalized Chi-square distribution with $`2L`$ degrees of freedom. Note that the fading power gains are exponentially distributed (Chi-square with 2 degrees of freedom) in the single transmit antenna case. We now characterize the asymptotic scaling laws of the throughput of the proposed static schedulers for this multi-transmit antenna scenario. Note that all the results in this section are derived for the case where $`L`$ is a constant and does not scale with $`N`$.
### 5.1 Worst User Scheduler
For the worst user scheme, the average throughput is given by
$$R_{tot}=NE\left[\mathrm{log}\left(1+|\chi _{min}|^2P\right)\right],$$
where $`|\chi _{min}|^2=\mathrm{min}_{i=1}^N|\chi _i|^2`$, and $`|\chi _i|^2`$ corresponds to the effective fading power gain at the $`i^{th}`$ user that follows a normalized Chi-square distribution with $`2L`$ degrees of freedom and whose distribution function is given by
$$F(x)=1e^{Lx}\left(\underset{k=0}{\overset{L1}{}}\frac{(Lx)^k}{k!}\right),x0.$$
(25)
###### Lemma 11
When the base station is equipped with $`L`$ transmit antennas, the average throughput of the worst user scheme scales as
$$R_{tot}=\mathrm{\Theta }\left(N^{\left(\frac{L1}{L}\right)}\right).$$
(26)
Thus the average throughput increases with $`L`$. This is expected since the performance of the worst user scheduler is degraded by the tail of the fading distribution. Hence, as $`L`$ increases, the spread of the fading distribution decreases, and consequently, the inherent multi-user diversity has a reduced effect on the performance of the scheduler. This leads to a rise in the average throughput of the worst user scheme from $`\mathrm{\Theta }(1)`$ for the single transmit antenna case to $`\mathrm{\Theta }(N)`$ for large values of $`L`$.
### 5.2 Best User Scheduler
For the best user scheme, the average throughput is given by
$$R_{tot}=E\left[\mathrm{log}\left(1+|\chi _{max}|^2P\right)\right],$$
where $`|\chi _{max}|^2=\mathrm{max}_{i=1}^N|\chi _i|^2`$.
###### Lemma 12
When the base station is equipped with $`L`$ transmit antennas, the average throughput of the best user scheme scales as
$$R_{tot}=\mathrm{\Theta }\left(\mathrm{log}\left(1+\frac{\mathrm{log}N+(L1)\mathrm{log}\mathrm{log}N}{L}\right)\right).$$
(27)
Since the best user scheduler leverages multi-user diversity to enhance the throughput, one can see that the throughput of the best user scheme decreases as $`L`$ increases.
## 6 Numerical Results
Here we present simulation results that validate our theoretical claims. These results were obtained through Monte-Carlo simulations and were averaged over at least 5000 iterations. The power constraint $`P`$ is taken to be unity. The throughput of the static schedulers, proposed in Section 3.1, is shown in Fig. 2 for different positions of the intended user in the ordered list of SNRs of all users. It is evident from the figure that, as predicted by the analysis, the throughput of the median user scheme is better than that of the best user scheme, which in turn outperforms the worst user scheme. In Fig. 3, we present a throughput-comparison for all the schemes proposed in Section 3 for increasing values of $`N`$. The corresponding delay-comparison is presented in Fig. 4. The throughput-comparison for the different scheduling schemes in the multi-group scenario is presented in Fig. 5 with $`G=5`$ groups (the corresponding delay-comparison is presented in Fig. 6). Although the best among worst users scheduler performs better than the best among best users scheme, in terms of throughput, for the range of $`N`$ values shown in the plot, it should be noted that the latter eventually outperforms the former for large values of $`N`$ ($`N>600`$). Except for this case, in all other considered scenarios, we can see that the simulation results follow the same trends predicted by our asymptotic analysis. Finally, we observe that the utility of our asymptotic analysis is manifested in its accurate predictions even with the relatively small number of users used in our simulations (i.e., in the order of $`N=10`$).
## 7 Conclusions
In this paper, we have used a cross layer design approach to shed more light on the throughput-delay tradeoff in the cellular multicast channel. Towards this end, we proposed three classes of scheduling algorithms with progressively increasing complexity, and analyzed the throughput-delay tradeoff achieved by each class. We first considered the class of low-complexity static scheduling schemes with memoryless decoding. We showed that a special case of this scheduling strategy, i.e., the median user scheduler, achieves the optimal scaling law of the throughput at the expense of an exponentially increasing delay with the number of users. We then proposed an incremental redundancy multicast scheme that achieves a superior throughput-delay tradeoff, at the expense of increased encoding/decoding complexity. We further proposed a cooperation scheme that achieves the optimal scaling laws of both throughput and delay at the expense of a high RF and computational complexity. We then generalized our schemes to the multi-group scenario and characterized their ability to exploit the multi-group diversity offered by the wireless channel. Finally, we presented simulation results that establish the accuracy of the predictions of our asymptotic analysis in systems with low to moderate number of users.
## Appendix A Proof of Theorem 3
The channel gain $`|h_{\pi (N\frac{N}{\alpha }+1)}|^2`$ has the distribution function $`F(x)`$ given by
$$F(x)=\underset{k=(N\frac{N}{\alpha }+1)}{\overset{N}{}}\left(\genfrac{}{}{0pt}{}{N}{k}\right)(1e^x)^ke^{(Nk)x},x0.$$
Hence the average throughput of the proposed scheme is given by
$$R_{tot}=\frac{N}{\alpha }_0^{\mathrm{}}\mathrm{log}(1+xP)dF(x).$$
Integrating by parts and simplifying, we obtain the average throughput as stated in equation (4) of the theorem.
We now calculate the average delay of the proposed scheme. We consider each coherence interval of length $`T_c`$ as a time slot. We first calculate the probability distribution of the service time $`X`$ required for transmitting a packet (of size $`S`$) when the base station always services the same queue. The service time $`X`$ is defined as
$$X=kT_c,k\{1,2,\mathrm{}\},$$
(28)
where $`k`$ is such that
$$T_c\left(\underset{i=1}{\overset{k1}{}}R_\alpha ^i\right)<ST_c\left(\underset{i=1}{\overset{k}{}}R_\alpha ^i\right).$$
(29)
Here $`R_\alpha ^i`$ represents the service rate in the $`i^{th}`$ time slot as given in (3). The probability distribution of $`X`$ is given by
$$\text{Pr}(X=kT_c)=\text{Pr}\left(\underset{i=1}{\overset{k1}{}}R_\alpha ^i<\frac{S}{T_c}\underset{i=1}{\overset{k}{}}R_\alpha ^i\right).$$
We let $`C=(S/T_c)`$ in the sequel. Using the exponential server assumption in (2) for the service rates $`\{R_\alpha ^i\}`$, we have (for $`k1`$)
$$\text{Pr}(X=kT_c)=_0^C\left(f_{\left(R_\alpha ^1+R_\alpha ^2+\mathrm{}+R_\alpha ^{(k1)}\right)}(y)dy\right)\text{Pr}(R_\alpha ^k>Cy)$$
$$=_0^C\frac{e^{\mu y}\mu ^{k1}y^{k2}}{(k2)!}e^{\mu (Cy)}๐y.$$
$$\text{Pr}(X=kT_c)=\frac{e^{\mu C}(\mu C)^{k1}}{(k1)!},k\{1,2,\mathrm{}\}.$$
(30)
Now the average service time $`\overline{X}`$ is given by
$$\overline{X}=\underset{k=1}{\overset{\mathrm{}}{}}kT_c\text{Pr}(X=kT_c)=\underset{k=1}{\overset{\mathrm{}}{}}kT_c\left(\frac{e^{\mu C}(\mu C)^{k1}}{(k1)!}\right).$$
$$\overline{X}=(1+\mu C)T_c=T_c+\mu S.$$
Since $`G=1`$ for the single group scenario, the assumption in (1) reduces to
$$T_c=\mathrm{\Theta }\left(\frac{1}{\mathrm{log}\mathrm{log}N}\right).$$
From the results in , we know that
$$E[R_\alpha ]E[R_N]=\mathrm{\Theta }(\mathrm{log}\mathrm{log}N).$$
Hence
$$\mu =\frac{1}{E[R_\alpha ]}=\mathrm{\Omega }\left(\frac{1}{\mathrm{log}\mathrm{log}N}\right),1\alpha N.$$
Thus for all possible values of the parameter $`\alpha `$, we have
$$\overline{X}=T_c+\mu S=\mathrm{\Theta }(\mu S)=\mathrm{\Theta }\left(\frac{1}{E[R_\alpha ]}\right).$$
Hence it is clear that the assumption on $`T_c`$ in (1) ensures that the average service time $`\overline{X}`$ is not dominated by the scaling behaviour of $`T_c`$.
We now focus on one set of $`\alpha `$ coupled queues. Any packet that arrives into this set enters all the $`\alpha `$ queues within the set and moreover, the base station services only one of the $`\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)`$ available queues in any time slot. Note that $`\overline{X}`$ was calculated assuming that the base station always services the same queue. We are interested in determining the delay involved in successfully transmitting a particular packet from all the $`\alpha `$ coupled queues in the set. The actual delay, as defined in Section 2, is the time between the start of transmission of a packet and the instant when the packet reaches all the $`N`$ users in the system. In our analysis, we assume that the packet of interest is at the head of all the $`\alpha `$ queues in the set during the start of transmission. This assumption thus yields a lower bound on the actual delay.
We characterize the delay based on the observation that our queuing problem is equivalent to the well-known โcoupon collectorโ problem. This observation was made earlier in where the authors characterized the delay of the throughput-optimal broadcast scheme. They assumed that the server (base station) offers a constant service rate which is independent of the instantaneous channel gains. In our analysis, however, we have incorporated the effects of rate adaptation. Let $`X_1,X_2,\mathrm{},X_\alpha `$ denote the service times (assuming continuous service), with distribution as given in (30), required for transmitting a packet from each of the $`\alpha `$ queues in the set. Then the delay of the proposed scheduling scheme is directly proportional to the minimum number of trials required to ensure that the first queue is served at least $`(X_1/T_c)`$ times by the base station, the second queue is served at least $`(X_2/T_c)`$ times and so on โฆ
We lower bound the average delay by calculating the minimum number of trials $`N_t`$ required to ensure that all the $`\alpha `$ queues are served at least $`(X_{min}/T_c)`$ times by the base station, where $`X_{min}=\mathrm{min}\{X_1,X_2,\mathrm{},X_\alpha \}`$. We determine the average number of such required trials $`E[N_t|X_{min}]`$ using the results derived in . Since the base station services only one of the $`\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)`$ queues in any time slot and the users are symmetric, there is an equal probability that the base station services any one of the queues. Since we need to consider only one set of $`\alpha `$ coupled queues for determining the delay, we consider all the other queues in the system jointly as one โdummyโ queue called the $`(\alpha +1)^{th}`$ queue. Now the probabilities $`\{p_j\}`$ of the server choosing the $`j^{th}`$ queue is given by
$$p_1=\mathrm{}=p_\alpha =\frac{1}{\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)}\text{and}p_{\alpha +1}=P_e=1\frac{\alpha }{\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)}.$$
These probabilities $`\{p_j\}`$ remain constant through all time slots and are not functions of the instantaneous service rates $`\{R_\alpha ^i\}`$ provided by the base station. The Moment Generating Function (MGF) of the number of trials required is given by
$$F_{N_t|X_{min}}(z)=\underset{i=0}{\overset{\mathrm{}}{}}z^i\text{Pr}(N_t>i)=\underset{i=0}{\overset{\mathrm{}}{}}z^ib_i,$$
where $`b_i`$ is the probability of failure of sending a packet to all the users in $`i`$ channel uses. The value of $`b_i`$ is equal to the polynomial
$$\left(\frac{x_1}{\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)}+\mathrm{}+\frac{x_\alpha }{\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)}+P_ex_{\alpha +1}\right)^i$$
evaluated at $`x_1=\mathrm{}=x_{\alpha +1}=1`$ after removing the terms that have all exponents of $`x_1,\mathrm{},x_\alpha `$ greater than or equal to $`(X_{min}/T_c)`$ (denoted by the operator $`\{.\}`$). Thus the MGF of the number of trials required is given by
$$F_{N_t|X_{min}}(z)=\underset{i=0}{\overset{\mathrm{}}{}}\frac{z^i}{\left[\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)\right]^i}\left\{\left(x_1+\mathrm{}+x_\alpha +P_e\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)x_{\alpha +1}\right)^i\right\}$$
evaluated at $`x_1=\mathrm{}=x_{\alpha +1}=1`$. But we know that
$$\frac{i!z^i}{\left[\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)\right]^i}=\frac{\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)}{z}_0^{\mathrm{}}e^{\frac{\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)t}{z}}t^i๐t.$$
Using this identity and simplifying, we get
$$F_{N_t|X_{min}}(z)=\frac{\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)}{z}_0^{\mathrm{}}e^{\frac{\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)t}{z}}\left(e^{\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)t}\left[1\left(1S_{\left(\frac{X_{min}}{T_c}\right)}(t)e^t\right)^\alpha \right]\right)๐t,$$
where
$$S_m(t)=\underset{i=0}{\overset{m1}{}}\frac{t^i}{i!}.$$
Hence the average number of trials required $`E[N_t|X_{min}]`$ is given by
$$E[N_t|X_{min}]=F_{N_t|X_{min}}(1)=\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)_0^{\mathrm{}}\left[1\left(1S_{\left(\frac{X_{min}}{T_c}\right)}(t)e^t\right)^\alpha \right]๐t=\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)E\left[\underset{1i\alpha }{\mathrm{max}}Y_i\right],$$
where the $`Y_i`$โs are i.i.d random variables that follow a Chi-square distribution with $`(2X_{min}/T_c)`$ degrees of freedom. From the results in , it can be seen that for such a sequence of random variables $`\{Y_i\}`$,
$$E\left[\underset{1i\alpha }{\mathrm{max}}Y_i\right]=\mathrm{max}\{\mathrm{\Theta }(\mathrm{log}\alpha ),\mathrm{\Theta }\left(\frac{X_{min}}{T_c}\right)\}.$$
(31)
Using this result, the average number of trials required is given by
$$E[N_t|X_{min}]=\mathrm{max}\{\mathrm{\Theta }\left(\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)\mathrm{log}\alpha \right),\mathrm{\Theta }\left(\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)\frac{X_{min}}{T_c}\right)\}.$$
Thus the average delay of the general static scheduling scheme can be lower bounded by
$$DE_{X_{min}}\left[E[N_t|X_{min}]T_c\right]=E_{X_{min}}\left[\mathrm{max}\{\mathrm{\Theta }\left(\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)T_c\mathrm{log}\alpha \right),\mathrm{\Theta }\left(\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)X_{min}\right)\}\right].$$
Since $`E\left[\mathrm{max}\{Z_1,Z_2\}\right]\mathrm{max}\{E[Z_1],E[Z_2]\}`$, we have
$$D=\mathrm{max}\{\mathrm{\Omega }\left(\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)T_c\mathrm{log}\alpha \right),\mathrm{\Omega }\left(\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)E[X_{min}]\right)\}.$$
$$D=\mathrm{max}\{\mathrm{\Omega }\left(\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)\frac{\mathrm{log}\alpha }{\mathrm{log}\mathrm{log}N}\right),\mathrm{\Omega }\left(\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)E[X_{min}]\right)\}.$$
(32)
Moreover, when $`E[X_{min}]=\mathrm{\Theta }\left(\overline{X}\right)`$, it can be easily seen that the expression on the right in (32) gives the exact scaling of the average delay $`D`$, instead of just being a lower bound for it.
## Appendix B Proof of Lemma 4
The average throughput of the worst user scheme is given by
$$R_{tot}=N_0^{\mathrm{}}\mathrm{log}(1+xP)Ne^{Nx}dx=Ne^{\left(\frac{N}{P}\right)}Ei\left(\frac{N}{P}\right).$$
(33)
For large values of $`x`$, we have
$$Ei(x)=_{\mathrm{}}^x\frac{e^t}{t}dt=\frac{e^x}{x}\left(1+ฯต\right),$$
where $`ฯต0`$ as $`x\mathrm{}`$. Using this in (33), we get
$$R_{tot}=P(1+ฯต)=\mathrm{\Theta }(1).$$
Letting $`\alpha =1`$ in (5) of Theorem 3, we get the average delay as<sup>5</sup><sup>5</sup>5Note that $`\mathrm{\Omega }\left(\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)\frac{\mathrm{log}\alpha }{\mathrm{log}\mathrm{log}N}\right)0`$ when $`\alpha =1`$, since for any constant $`k`$, $`k+\mathrm{log}\alpha =\mathrm{\Theta }(\mathrm{log}\alpha )`$ $`1\alpha N`$.
$$D=\mathrm{max}\{\mathrm{\Omega }\left(\frac{1}{\mathrm{log}\mathrm{log}N}\right),\mathrm{\Omega }\left(E[X_{min}]\right)\}.$$
(34)
Since the base station maintains only a single queue for the worst user scheme, we have
$$E[X_{min}]=\overline{X}=\mathrm{\Theta }\left(\frac{1}{E[R_1]}\right).$$
The average service rate $`E[R_1]`$ is given by
$$E[R_1]=_0^{\mathrm{}}\mathrm{log}(1+xP)Ne^{Nx}dx=e^{\left(\frac{N}{P}\right)}Ei\left(\frac{N}{P}\right)=\mathrm{\Theta }\left(\frac{1}{N}\right).$$
Since $`E[X_{min}]=\overline{X}`$, the expression on the right in (34) gives the exact scaling of $`D`$. Thus the average delay of the worst user scheme scales as $`D=\mathrm{\Theta }(N)`$.
## Appendix C Proof of Lemma 5
By letting $`\alpha =N`$ in (4) of Theorem 3, the average throughput of the best user scheme is found to be
$$R_{tot}=\underset{i=1}{\overset{N}{}}\left(\genfrac{}{}{0pt}{}{N}{i}\right)(1)^ie^{\left(\frac{i}{P}\right)}Ei\left(\frac{i}{P}\right).$$
(35)
It has been shown in that the throughput in (35) scales as
$$R_{tot}=\mathrm{\Theta }(\mathrm{log}\mathrm{log}N)$$
(36)
with the number of users $`N`$. Hence the average service rate is given by
$$E[R_N]=R_{tot}=\mathrm{\Theta }\left(\mathrm{log}\mathrm{log}N\right).$$
Letting $`\alpha =N`$ in (5) of Theorem 3, we get the average delay as
$$D=\mathrm{max}\{\mathrm{\Omega }\left(\frac{N\mathrm{log}N}{\mathrm{log}\mathrm{log}N}\right),\mathrm{\Omega }\left(NE[X_{min}]\right)\}.$$
(37)
Now
$$E[X_{min}]=E\left[\underset{i=1}{\overset{N}{\mathrm{min}}}X_i\right]\overline{X}=\mathrm{\Theta }\left(\frac{1}{E[R_N]}\right).$$
Hence
$$E[X_{min}]=O\left(\frac{1}{\mathrm{log}\mathrm{log}N}\right).$$
Thus, from (37), the average delay of the best user scheme scales as
$$D=\mathrm{\Omega }\left(\frac{N\mathrm{log}N}{\mathrm{log}\mathrm{log}N}\right).$$
## Appendix D Proof of Lemma 6
The average throughput of the median user scheme can be derived by letting $`\alpha =2`$ in (4) of Theorem 3. From the results on central order statistics in (Theorem 8.5.1), we know that the sample median of $`N`$ i.i.d. exponential random variables converges in distribution to a normal random variable with mean $`\theta `$ and variance $`(1/N)`$, where $`\theta =\mathrm{log}2`$ is the median of the underlying exponential distribution. Hence
$$\left(|h_{\pi (\frac{N}{2}+1)}|^2\theta \right)\sqrt{N}W\text{in distribution},$$
where $`W`$ is a standard normal random variable. Using Chebyshevโs inequality, we get $`ฯต>0`$
$$\text{Pr}\left(\left||h_{\pi (\frac{N}{2}+1)}|^2\theta \right|>ฯต\right)=\text{Pr}\left(\sqrt{N}\left||h_{\pi (\frac{N}{2}+1)}|^2\theta \right|>ฯต\sqrt{N}\right)<\frac{E[W^2]+\delta }{Nฯต^2}0\text{as }N\mathrm{}.$$
$$|h_{\pi (\frac{N}{2}+1)}|^2\theta \text{in probability}.$$
Since the $`\mathrm{log}(.)`$ function is continuous, we have
$$\mathrm{log}\left(1+|h_{\pi (\frac{N}{2}+1)}|^2P\right)\mathrm{log}(1+\theta P)\text{in probability}.$$
(38)
We now derive a lower bound on the average throughput. We recall the following property of positive random variables. Let $`(X_n)`$ be a set of positive random variables converging to a constant $`A`$ in probability. Hence $`ฯต>0`$,
$$\text{Pr}\left(|X_nA|ฯต\right)<\delta ,$$
for some small $`\delta >0`$. Now
$$E[X_n]=_0^{\mathrm{}}tf_{X_n}(t)\mathrm{d}t_{Aฯต}^{A+ฯต}tf_{X_n}(t)\mathrm{d}t(Aฯต)(1\delta ).$$
Taking the limit as $`n\mathrm{}`$, we get
$$\underset{n\mathrm{}}{lim}E[X_n]A.$$
(39)
Using this property in (38), we get
$$\underset{N\mathrm{}}{lim}E\left[\mathrm{log}\left(1+|h_{\pi (\frac{N}{2}+1)}|^2P\right)\right]\mathrm{log}(1+\theta P)=\mathrm{\Theta }(1).$$
$$R_{tot}=\left(\frac{N}{2}\right)E\left[\mathrm{log}\left(1+|h_{\pi (\frac{N}{2}+1)}|^2P\right)\right]=\mathrm{\Omega }(N).$$
(40)
An upper bound on the average throughput of any scheduling scheme is given by
$$R_{tot}E\left[\underset{i=1}{\overset{N}{}}\mathrm{log}\left(1+|h_i|^2P\right)\right]=NE\left[\mathrm{log}\left(1+|h_1|^2P\right)\right].$$
Hence
$$R_{tot}=O(N).$$
(41)
Combining this with the lower bound in (40), we get
$$R_{tot}=\mathrm{\Theta }(N).$$
Thus it is clear that the throughput of the proposed median user scheme is scaling law optimal. Letting $`\alpha =2`$ in (5) of Theorem 3, we get the average delay as
$$D=\mathrm{max}\{\mathrm{\Omega }\left(\left(\genfrac{}{}{0pt}{}{N}{N/2}\right)\frac{1}{\mathrm{log}\mathrm{log}N}\right),\mathrm{\Omega }\left(\left(\genfrac{}{}{0pt}{}{N}{N/2}\right)E[X_{min}]\right)\}.$$
(42)
Now, since $`X_{min}=\mathrm{min}\{X_1,X_2\}`$, we have
$$E[X_{min}]=\mathrm{\Theta }\left(\overline{X}\right)=\mathrm{\Theta }\left(\frac{1}{E[R_2]}\right).$$
The average service rate $`E[R_2]`$ is given by
$$E[R_2]=E\left[\mathrm{log}\left(1+|h_{\pi (\frac{N}{2}+1)}|^2P\right)\right]=\mathrm{\Theta }(1).$$
(43)
Since $`E[X_{min}]=\mathrm{\Theta }\left(\overline{X}\right)`$, the expression on the right in (42) gives the exact scaling of $`D`$. Thus the average delay of the median user scheme scales as
$$D=\mathrm{\Theta }\left(\left(\genfrac{}{}{0pt}{}{N}{N/2}\right)\right).$$
Using Stirlingโs formula, we obtain the scaling of the average delay as given in (11).
## Appendix E Proof of Theorem 7
Let $`A_i`$ denote the event that a packet is successfully decoded by all the $`N`$ users in the system in $`i`$ transmission attempts. Following the notation in , we define
$$q(m)=\text{Pr}(\overline{A_1},\mathrm{},\overline{A_{m1}},A_m)=p(m1)p(m),$$
where
$$p(m)=\text{Pr}(\overline{A_1},\mathrm{},\overline{A_{m1}},\overline{A_m})=1\underset{l=1}{\overset{m}{}}q(l),$$
with $`p(0)=1`$. The rate $`\overline{R}`$ is defined as $`\overline{R}=(b/L)`$. We define the random variable $`\tau `$ to be the number of transmission attempts made between the instant when the codeword is generated and the instant when its transmission is stopped (Transmission is stopped either when the packet is successfully decoded by all the $`N`$ users or the number of transmission attempts exceeds the rate constraint $`M`$). The probability distribution of $`\tau `$ is given by
$$f_\tau (m)=\{\begin{array}{c}0,m=0\hfill \\ q(m),1mM1\hfill \\ q(M)+p(M),m=M\hfill \end{array}$$
We define the random reward $``$ as follows: $`=N\overline{R}`$ if transmission stops because of successful decoding and $`=0`$ if transmission stops because of the rate constraint violation. Hence
$$E[]=N\overline{R}\underset{m=1}{\overset{M}{}}q(m)=N\overline{R}[1p(M)].$$
The mean inter-renewal time is given by
$$E[\tau ]=\underset{m=1}{\overset{M}{}}mf_\tau (m)=\underset{m=1}{\overset{M}{}}mq(m)+Mp(M)=\underset{m=1}{\overset{M}{}}m[p(m1)p(m)]+Mp(M)=\underset{m=0}{\overset{M1}{}}p(m).$$
Applying the renewal-reward theorem, we obtain the average throughput of the proposed scheme as
$$R_{tot}=\frac{E[]}{E[\tau ]},\text{with probability }1.$$
Hence
$$R_{tot}=\frac{N\overline{R}\left[1p(M)\right]}{1+_{m=1}^{M1}p(m)}.$$
The average delay $`D`$ of the scheme is given by the mean inter-renewal time. Hence
$$D=E[\tau ]=1+\underset{m=1}{\overset{M1}{}}p(m).$$
The unconstrained throughput and delay of this scheme are obtained by letting $`M\mathrm{}`$ and are given by
$$R_{tot}=\frac{N\overline{R}}{_{m=0}^{\mathrm{}}p(m)},$$
(44)
and
$$D=\underset{m=0}{\overset{\mathrm{}}{}}p(m).$$
(45)
From the earlier definitions, we have
$$p(m)=\text{Pr}(\overline{A_1},\mathrm{},\overline{A_{m1}},\overline{A_m})=\text{Pr}(\overline{A_m})=\text{Pr}\left(\underset{i=1}{\overset{N}{\mathrm{min}}}\underset{k=1}{\overset{m}{}}I(X;Y_{ik})\overline{R}\right).$$
$$p(m)=1\text{Pr}\left(\underset{i=1}{\overset{N}{\mathrm{min}}}\underset{k=1}{\overset{m}{}}I(X;Y_{ik})>\overline{R}\right)=1\left[1\text{Pr}\left(\underset{k=1}{\overset{m}{}}I(X;Y_{1k})\overline{R}\right)\right]^N.$$
(46)
Now for a Gaussian input distribution, we have
$$\underset{k=1}{\overset{m}{}}I(X;Y_{1k})=\underset{k=1}{\overset{m}{}}\mathrm{log}(1+|h_k|^2).$$
We know that
$$\mathrm{log}\left(1+\underset{k=1}{\overset{m}{}}|h_k|^2\right)\underset{k=1}{\overset{m}{}}\mathrm{log}(1+|h_k|^2)\underset{k=1}{\overset{m}{}}|h_k|^2.$$
Hence
$$\text{Pr}\left(\underset{k=1}{\overset{m}{}}|h_k|^2(e^{\overline{R}}1)\right)\text{Pr}\left(\underset{k=1}{\overset{m}{}}\mathrm{log}(1+|h_k|^2)\overline{R}\right)\text{Pr}\left(\underset{k=1}{\overset{m}{}}|h_k|^2\overline{R}\right).$$
Since both $`\overline{R}`$ and $`(e^{\overline{R}}1)`$ are constants, substituting both the lower and upper bounds in (46) will yield the same scaling with $`N`$. So we consider only the lower bound on $`p(m)`$. Let
$$s(m)=1\left[1\text{Pr}\left(\underset{k=1}{\overset{m}{}}|h_k|^2\overline{R}\right)\right]^N.$$
Hence
$$\underset{m=0}{\overset{\mathrm{}}{}}p(m)=\mathrm{\Theta }\left(\underset{m=0}{\overset{\mathrm{}}{}}s(m)\right)\text{w.r.t }N.$$
The random variable $`_{k=1}^m|h_k|^2`$ has a $`2m`$-dimensional Chi-square distribution with the density and distribution functions given by
$$f(x)=\frac{e^xx^{m1}}{(m1)!},x0$$
and
$$F(x)=1e^x\left(\underset{l=0}{\overset{m1}{}}\frac{x^l}{l!}\right),x0.$$
Hence
$$s(m)=1\left[e^{\overline{R}}\left(\underset{l=0}{\overset{m1}{}}\frac{\overline{R}^l}{l!}\right)\right]^N.$$
From Taylorโs theorem, we know that (for some $`0<\theta <1`$)
$$e^{\overline{R}}=\underset{l=0}{\overset{m1}{}}\frac{\overline{R}^l}{l!}+\frac{e^{\theta \overline{R}}\overline{R}^m}{m!}\underset{l=0}{\overset{m1}{}}\frac{\overline{R}^l}{l!}=e^{\overline{R}}\frac{e^{\theta \overline{R}}\overline{R}^m}{m!}.$$
$$s(m)=1\left(1\frac{e^{(1\theta )\overline{R}}\overline{R}^m}{m!}\right)^N.$$
To find the scaling of $`_{m=0}^{\mathrm{}}s(m)`$ w.r.t $`N`$, we first derive a lower bound by finding the value of $`m`$ until which $`s(m)1`$ as $`N\mathrm{}`$. Now
$$s(m)1(1\frac{e^{(1\theta )\overline{R}}\overline{R}^m}{m!})^N0.$$
$$\frac{e^{(1\theta )\overline{R}}\overline{R}^m}{m!}>\mathrm{\Theta }\left(\frac{1}{N}\right).$$
Using Stirlingโs approximation, we have
$$\frac{e^{(1\theta )\overline{R}}\overline{R}^m}{\sqrt{2\pi m}e^mm^m}>\frac{k}{N},\text{constant }k.$$
Taking log on both sides, we get
$$(1\theta )\overline{R}m+m\mathrm{log}\left(\frac{m}{\overline{R}}\right)+\frac{1}{2}\mathrm{log}(2\pi m)<\mathrm{log}N\mathrm{log}k,k.$$
For large $`N`$, this equation can be reduced to
$$m\mathrm{log}m<\mathrm{log}N.$$
(47)
This equation is satisfied by all values of $`m`$ such that
$$m<\mathrm{\Theta }\left(\frac{\mathrm{log}N}{\mathrm{log}\mathrm{log}N}\right).$$
Since $`s(m)1`$ as $`N\mathrm{}`$ for all values of $`m`$ that satisfy the above equation, the sum of $`s(m)`$โs can be lower bounded as
$$\underset{m=0}{\overset{\mathrm{}}{}}s(m)\mathrm{\Theta }\left(\frac{\mathrm{log}N}{\mathrm{log}\mathrm{log}N}\right).$$
(48)
Similarly an upper bound on $`_{m=0}^{\mathrm{}}s(m)`$ can be derived by finding the value of $`m`$ from which $`s(m)0`$ as $`N\mathrm{}`$. Following the same procedure as before, we find that $`s(m)0`$ when
$$m>\mathrm{\Theta }\left(\frac{\mathrm{log}N}{\mathrm{log}\mathrm{log}N}\right).$$
This yields the following upper bound
$$\underset{m=0}{\overset{\mathrm{}}{}}s(m)\mathrm{\Theta }\left(\frac{\mathrm{log}N}{\mathrm{log}\mathrm{log}N}\right).$$
Combining this with the lower bound in (48), we get
$$\underset{m=0}{\overset{\mathrm{}}{}}s(m)=\mathrm{\Theta }\left(\frac{\mathrm{log}N}{\mathrm{log}\mathrm{log}N}\right).$$
Thus the average delay is given by
$$D=\underset{m=0}{\overset{\mathrm{}}{}}p(m)=\mathrm{\Theta }\left(\underset{m=0}{\overset{\mathrm{}}{}}s(m)\right)=\mathrm{\Theta }\left(\frac{\mathrm{log}N}{\mathrm{log}\mathrm{log}N}\right).$$
The average throughput of the incremental redundancy scheme is then given by
$$R_{tot}=\frac{N\overline{R}}{_{m=0}^{\mathrm{}}p(m)}=\frac{N\overline{R}}{D}=\mathrm{\Theta }\left(\frac{N\mathrm{log}\mathrm{log}N}{\mathrm{log}N}\right).$$
## Appendix F Proof of Theorem 8
The first stage of the cooperation scheme is the median user scheme. Hence it is clear from (10) that
$$E[R_{s1}]=E\left[\mathrm{log}\left(1+|h_{\pi (\frac{N}{2}+1)}|^2P\right)\right]=\mathrm{\Theta }(1).$$
As noted earlier, the cooperative transmission by the users in the second stage is equivalent to the transmission of packets from a transmitter equipped with $`(N/2)`$ transmit antennas to the worst user in a group of $`(N/2)`$ users. Hence the average transmission rate during the cooperative stage is given by
$$E[R_{s2}]=E\left[\underset{i=1,\mathrm{},(N/2)}{\mathrm{min}}\mathrm{log}\left(1+\frac{|h_{1i}|^2+\mathrm{}+|h_{(N/2)i}|^2}{(N/2)}P\right)\right],$$
where the $`|h_{ki}|^2`$โs are i.i.d and exponentially distributed and represent the inter-user fading coefficients.
$$E[R_{s2}]=E\left[\mathrm{log}\left(1+\underset{i=1,\mathrm{},M}{\mathrm{min}}\frac{|\chi _{2M}^i|^2}{M}P\right)\right],$$
(49)
where $`M=(N/2)`$ and $`|\chi _{2M}^i|^2`$โs are Chi-square random variables with $`2M`$ degrees of freedom whose distribution function is given by
$$F(x)=1e^x\left(\underset{j=0}{\overset{M1}{}}\frac{x^j}{j!}\right),x0.$$
Using the results on extreme order statistics in (Theorems 8.3.2-8.3.6), it can be shown that the random variable
$$\frac{\mathrm{min}_{i=1}^M|\chi _{2M}^i|^2}{b_M}W\text{in distribution as }M\mathrm{},$$
where $`W`$ is a Weibull type random variable and $`b_M`$ satisfies $`F(b_M)=\frac{1}{M}`$. Now
$$F(b_M)=\frac{1}{M}1e^{b_M}\left(\underset{j=0}{\overset{M1}{}}\frac{b_M^j}{j!}\right)=\frac{1}{M}.$$
Using Taylorโs theorem, we get for some $`0<\beta _M<1`$
$$1e^{b_M}(e^{b_M}\frac{e^{\beta _Mb_M}b_M^M}{M!})=\frac{1}{M}\frac{e^{(1\beta _M)b_M}b_M^M}{M!}=\frac{1}{M}.$$
Using Stirlingโs approximation, we have
$$\frac{e^{(1\beta _M)b_M}b_M^M}{\sqrt{2\pi M}M^Me^M}=\frac{1}{M}.$$
Taking $`\mathrm{log}(.)`$ on both sides, we get
$$(1\beta _M)b_MM\mathrm{log}b_M=M\left(M\frac{1}{2}\right)\mathrm{log}M+C.$$
Since $`\beta _M0`$ as $`M\mathrm{}`$, we get $`b_M=\mathrm{\Theta }(M)`$. Thus
$$\frac{\mathrm{min}_{i=1}^M|\chi _{2M}^i|^2}{M}kW\text{in distribution, for some constant }k>0.$$
Since the $`\mathrm{log}(.)`$ function is continuous, we have
$$\mathrm{log}\left(1+\frac{\mathrm{min}_{i=1}^M|\chi _{2M}^i|^2}{M}P\right)\mathrm{log}(1+kWP)\text{in distribution, as }M\mathrm{}.$$
Now, we know
$$\mathrm{log}\left(1+\frac{\mathrm{min}_{i=1}^M|\chi _{2M}^i|^2}{M}P\right)\frac{\left(\mathrm{min}_{i=1}^M|\chi _{2M}^i|^2\right)P}{M}\frac{|\chi _{2M}^1|^2P}{M}.$$
Since
$$E\left[\left(\frac{|\chi _{2M}^1|^2P}{M}\right)^2\right]=\frac{E[(|\chi _{2M}^1|^2)^2]P^2}{M^2}=\left(\frac{M^2+M}{M^2}\right)P^2=\left(1+\frac{1}{M}\right)P^22P^2<\mathrm{}M,$$
$$\left\{\frac{|\chi _{2M}^1|^2P}{M};M1\right\}\text{is uniformly integrable.}$$
$$\left\{\mathrm{log}\left(1+\frac{\mathrm{min}_{i=1}^M|\chi _{2M}^i|^2}{M}P\right);M1\right\}\text{is uniformly integrable.}$$
It is shown in that if a sequence of random variables $`(X_n)`$ is uniformly integrable and $`X_nX`$ in distribution as $`n\mathrm{}`$, then $`EX_nEX`$ as $`n\mathrm{}`$. Thus
$$E\left[\mathrm{log}\left(1+\frac{\mathrm{min}_{i=1}^M|\chi _{2M}^i|^2}{M}P\right)\right]E[\mathrm{log}(1+kWP)]=\mathrm{\Theta }(1).$$
Hence the average transmission rate of the second stage is given by $`E[R_{s2}]=\mathrm{\Theta }(1)`$ w.r.t $`N`$. Since both $`E[R_{s1}]`$ and $`E[R_{s2}]`$ do not scale with $`N`$ and since the minimum is taken over only two positive quantities, we have
$$E\left[\mathrm{min}\{R_{s1},R_{s2}\}\right]=\mathrm{\Theta }\left(E[R_{s1}]\right)=\mathrm{\Theta }(1).$$
Thus the average throughput of the cooperation scheme is given by
$$R_{tot}=\left(\frac{N}{2}\right)E\left[\mathrm{min}\{R_{s1},R_{s2}\}\right]=\mathrm{\Theta }(N).$$
We now determine the average delay of the cooperation scheme. We note that the base station needs to maintain only a single queue that caters to all the $`N`$ users in the system. The information transmitted by the base station in the first half of each time slot reaches all the $`N`$ users at the end of that time slot. Hence the average delay is equal to the average service time required for transmitting a packet of size $`S`$ from the queue. Following the steps in Appendix A, the average delay $`D`$ for transmitting a packet in the cooperation scheme is given by
$$D=T_c+\mu S=\mathrm{\Theta }\left(\frac{S}{E[\mathrm{min}\{R_{s1},R_{s2}\}]}\right)=\mathrm{\Theta }(1).$$
## Appendix G Proof of Theorem 9
We first extend the proof of the general static scheduling scheme given in Appendix A to the multi-group scenario and then consider the three special cases. The average throughput $`R_{tot}`$ of the general multi-group scheduling scheme is given by
$$R_{tot}=\left(\frac{N}{\alpha }\right)E[R_\alpha ],$$
where $`R_\alpha `$ represents the transmission rate to each of the intended $`(N/\alpha )`$ users and is given by
$$R_\alpha =\mathrm{log}\left(1+|h_{bg}|^2P\right),$$
where the distribution of $`|h_{bg}|^2`$ is given by
$$F_{bg}(x)=\left(\underset{k=(N\frac{N}{\alpha }+1)}{\overset{N}{}}\left(\genfrac{}{}{0pt}{}{N}{k}\right)\left(1e^x\right)^ke^{(Nk)x}\right)^G,x0.$$
(50)
Hence the average throughput is given by
$$R_{tot}=\frac{N}{\alpha }_0^{\mathrm{}}\mathrm{log}(1+xP)dF_{bg}(x).$$
(51)
Integrating by parts and simplifying, we obtain the average throughput of the proposed scheme.
For implementing the general multi-group static scheduling scheme, the base station needs to maintain $`G\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)`$ queues, one for each combination of $`(N/\alpha )`$ users in each of the $`G`$ groups. These queues can be divided into $`G`$ sets, one for each of the groups. Within each set corresponding to a particular group, the queues can be further divided into subsets with $`\alpha `$ coupled queues in each subset such that the combinations of users served by the $`\alpha `$ queues within a subset are mutually exclusive and collectively exhaustive (i.e., every user in the particular group is served by exactly one of the $`\alpha `$ queues). We consider one such subset of $`\alpha `$ queues corresponding to any one of the $`G`$ groups. Any packet that arrives into the subset enters all the $`\alpha `$ queues since it needs to be transmitted to all the users within the group. At any instant of time, the base station services only one of the $`G\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)`$ queues.
As before, we first calculate the average service time $`\overline{X}`$ required for transmitting a packet by assuming that the base station always services the same queue. Following the steps in Appendix A, the average service time $`\overline{X}`$ is given by
$$\overline{X}=T_c+\mu S=\mathrm{\Theta }\left(\frac{S}{E[R_\alpha ]}\right).$$
We again use the results in to derive a lower bound on the actual delay by considering the minimum number of trials $`N_t`$ required to ensure that all the $`\alpha `$ queues are served at least $`(X_{min}/T_c)`$ times by the base station. As before, we consider only $`(\alpha +1)`$ queues with the $`(\alpha +1)^{th}`$ queue being the โdummyโ queue representing all the queues in all other subsets in the system. Now the probabilities $`\{p_i\}`$ of the server choosing the $`i^{th}`$ queue are given by
$$p_1=\mathrm{}=p_\alpha =\frac{1}{G\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)}\text{and}p_{\alpha +1}=1\frac{\alpha }{G\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)}.$$
Proceeding as in Appendix A, the MGF of the number of trials required is given by
$$F_{N_t|X_{min}}(z)=\frac{G\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)}{z}_0^{\mathrm{}}e^{\frac{G\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)t}{z}}\left(e^{G\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)t}\left[1\left(1S_{\left(\frac{X_{min}}{T_c}\right)}(t)e^t\right)^\alpha \right]\right)๐t.$$
The average number of trials required $`E[N_t|X_{min}]`$ is given by
$$E[N_t|X_{min}]=F_{N_t|X_{min}}(1)=G\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)_0^{\mathrm{}}\left[1\left(1S_{\left(\frac{X_{min}}{T_c}\right)}(t)e^t\right)^\alpha \right]๐t=G\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)E\left[\underset{1i\alpha }{\mathrm{max}}Y_i\right],$$
where the $`Y_i`$โs are i.i.d random variables that follow a Chi-square distribution with $`(2X_{min}/T_c)`$ degrees of freedom. Using the result in (31), we get
$$E[N_t|X_{min}]=\mathrm{max}\{\mathrm{\Theta }\left(G\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)\mathrm{log}\alpha \right),\mathrm{\Theta }\left(G\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)\frac{X_{min}}{T_c}\right)\}.$$
Thus the average delay of the proposed scheme can be lower bounded as
$$DE_{X_{min}}\left[E[N_t|X_{min}]T_c\right]=E_{X_{min}}\left[\mathrm{max}\{\mathrm{\Theta }\left(G\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)T_c\mathrm{log}\alpha \right),\mathrm{\Theta }\left(G\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)X_{min}\right)\}\right].$$
$$D=\mathrm{max}\{\mathrm{\Omega }\left(\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)\frac{G\mathrm{log}\alpha }{\mathrm{log}\mathrm{log}NG}\right),\mathrm{\Omega }\left(G\left(\genfrac{}{}{0pt}{}{N}{N/\alpha }\right)E[X_{min}]\right)\}.$$
(52)
Moreover, when $`E[X_{min}]=\mathrm{\Theta }\left(\overline{X}\right)`$, the expression on the right in (52) gives the exact scaling of the average delay $`D`$, instead of just being a lower bound for it.
### G.1 Best among worst users scheme ($`\alpha =1`$)
Letting $`\alpha =1`$ in (51) and simplifying, we get the average throughput to be
$$R_{tot}=N\left[\underset{k=1}{\overset{G}{}}\left(\genfrac{}{}{0pt}{}{G}{k}\right)(1)^ke^{\left(\frac{Nk}{P}\right)}Ei\left(\frac{Nk}{P}\right)\right].$$
(53)
For large values of $`x`$, we have
$$Ei(x)=_{\mathrm{}}^x\frac{e^t}{t}dt=\frac{e^x}{x}\left(1+ฯต\right),$$
where $`ฯต0`$ as $`x\mathrm{}`$. Using this in (53), we get
$$R_{tot}=P\left[\underset{k=1}{\overset{G}{}}\left(\genfrac{}{}{0pt}{}{G}{k}\right)\frac{(1)^{k+1}}{k}\right](1+ฯต).$$
It can be shown using the results in that
$$R_{tot}=P\left[\underset{k=1}{\overset{G}{}}\frac{1}{k}\right](1+ฯต)=\mathrm{\Theta }(\mathrm{log}G).$$
The average service rate is given by
$$E[R_1]=\underset{k=1}{\overset{G}{}}\left(\genfrac{}{}{0pt}{}{G}{k}\right)(1)^ke^{\left(\frac{Nk}{P}\right)}Ei\left(\frac{Nk}{P}\right)=\mathrm{\Theta }\left(\frac{\mathrm{log}G}{N}\right).$$
Letting $`\alpha =1`$ in (52) and using the fact that
$$E[X_{min}]=\overline{X}=\mathrm{\Theta }\left(\frac{1}{E[R_1]}\right),$$
we get the average delay as
$$D=\mathrm{\Theta }\left(\frac{NG}{\mathrm{log}G}\right).$$
### G.2 Best among best users scheme ($`\alpha =N`$)
Letting $`\alpha =N`$ in (51) and simplifying, we get the average throughput to be
$$R_{tot}=\underset{k=1}{\overset{NG}{}}\left(\genfrac{}{}{0pt}{}{NG}{k}\right)(1)^ke^{\left(\frac{k}{P}\right)}Ei\left(\frac{k}{P}\right).$$
(54)
It has been shown in that the throughput in (54) scales as
$$R_{tot}=\mathrm{\Theta }(\mathrm{log}\mathrm{log}NG).$$
(55)
Hence the average service rate is given by
$$E[R_N]=R_{tot}=\mathrm{\Theta }\left(\mathrm{log}\mathrm{log}NG\right).$$
Letting $`\alpha =N`$ in (52) and using the fact that
$$E[X_{min}]\overline{X}=\mathrm{\Theta }\left(\frac{1}{E[R_N]}\right),$$
we get the average delay as
$$D=\mathrm{\Omega }\left(\frac{NG\mathrm{log}N}{\mathrm{log}\mathrm{log}NG}\right).$$
### G.3 Best among median users scheme ($`\alpha =2`$)
The average throughput of the best among median users scheme can be derived by letting $`\alpha =2`$ in (51). It is given by
$$R_{tot}=\left(\frac{N}{2}\right)E\left[\mathrm{log}\left(1+\underset{g=1}{\overset{G}{\mathrm{max}}}|h_{\pi (\frac{N}{2}+1)}^g|^2P\right)\right].$$
We now determine bounds on the asymptotic scaling of $`R_{tot}`$ as $`N`$ and $`G`$ grow to infinity. To get a lower bound on the throughput, we use the fact that $`\mathrm{max}\{X_1,\mathrm{},X_n\}X_1`$ and obtain
$$R_{tot}\left(\frac{N}{2}\right)E\left[\mathrm{log}\left(1+|h_{\pi (\frac{N}{2}+1)}^1|^2P\right)\right].$$
The expression on the right is clearly the throughput of the median user scheduler described in Section 3.1. Hence, from Lemma 6, we get
$$R_{tot}=\mathrm{\Omega }(N).$$
We derive an upper bound on the throughput by using the fact that for continuous unimodal distributions
$$\left|\text{Median}\text{Mean}\right|<\text{Standard Deviation}.$$
Hence<sup>6</sup><sup>6</sup>6It is easy to show using convergence arguments that the inequality is valid for the empirical values used in the proof.
$$|h_{\pi (\frac{N}{2}+1)}|^2<\frac{|h_1|^2+\mathrm{}+|h_N|^2}{N}+\sqrt{\left(\frac{|h_1|^4+\mathrm{}+|h_N|^4}{N}\right)\left(\frac{|h_1|^2+\mathrm{}+|h_N|^2}{N}\right)^2}.$$
$$E\left[\underset{g=1}{\overset{G}{\mathrm{max}}}|h_{\pi (\frac{N}{2}+1)}^g|^2\right]<E\left[\underset{g=1}{\overset{G}{\mathrm{max}}}\left\{\frac{|h_1^g|^2+\mathrm{}+|h_N^g|^2}{N}+\sqrt{\frac{|h_1^g|^4+\mathrm{}+|h_N^g|^4}{N}}\right\}\right].$$
Using Jensenโs inequality and the fact that $`\mathrm{max}_{i=1}^n\{X_i+Y_i\}\mathrm{max}_{i=1}^n\{X_i\}+max_{i=1}^n\{Y_i\}`$, we get
$$E\left[\underset{g=1}{\overset{G}{\mathrm{max}}}|h_{\pi (\frac{N}{2}+1)}^g|^2\right]<E\left[\underset{g=1}{\overset{G}{\mathrm{max}}}\left\{\frac{|h_1^g|^2+\mathrm{}+|h_N^g|^2}{N}\right\}\right]+\sqrt{E\left[\underset{g=1}{\overset{G}{\mathrm{max}}}\left\{\frac{|h_1^g|^4+\mathrm{}+|h_N^g|^4}{N}\right\}\right]}$$
$$<\frac{1}{N}\underset{i=1}{\overset{N}{}}E\left[\underset{g=1}{\overset{G}{\mathrm{max}}}|h_i^g|^2\right]+\sqrt{\frac{1}{N}\underset{i=1}{\overset{N}{}}E\left[\underset{g=1}{\overset{G}{\mathrm{max}}}|h_i^g|^4\right]}=E\left[\underset{g=1}{\overset{G}{\mathrm{max}}}|h_1^g|^2\right]+\sqrt{E\left[\left(\underset{g=1}{\overset{G}{\mathrm{max}}}|h_1^g|^2\right)^2\right]}.$$
It is known that for a sequence of exponential random variables $`\{X_i\}`$ with unit mean ,
$$E\left[\underset{i=1}{\overset{G}{\mathrm{max}}}X_i\right]=\mathrm{\Theta }(\mathrm{log}G)\text{and}E\left[\left(\underset{i=1}{\overset{G}{\mathrm{max}}}X_i\right)^2\right]=\mathrm{\Theta }\left((\mathrm{log}G)^2\right).$$
Thus
$$E\left[\underset{g=1}{\overset{G}{\mathrm{max}}}|h_{\pi (\frac{N}{2}+1)}^g|^2\right]=O(\mathrm{log}G).$$
Now applying Jensenโs inequality, we get the upper bound on the average throughput as
$$R_{tot}=\left(\frac{N}{2}\right)E\left[\mathrm{log}\left(1+\underset{g=1}{\overset{G}{\mathrm{max}}}|h_{\pi (\frac{N}{2}+1)}^g|^2P\right)\right]\left(\frac{N}{2}\right)\mathrm{log}\left(1+E\left[\underset{g=1}{\overset{G}{\mathrm{max}}}|h_{\pi (\frac{N}{2}+1)}^g|^2\right]P\right).$$
$$R_{tot}=O(N\mathrm{log}\mathrm{log}G).$$
Thus the average throughput of the best among median users scheme can be bounded as
$$\mathrm{\Omega }(N)=R_{tot}=O(N\mathrm{log}\mathrm{log}G).$$
(56)
The average service rate $`E[R_2]`$ is then bounded by
$$\mathrm{\Omega }(1)=E[R_2]=O(\mathrm{log}\mathrm{log}G).$$
Letting $`\alpha =2`$ in (52) and using the fact that
$$E[X_{min}]=\mathrm{\Theta }\left(\overline{X}\right)=\mathrm{\Theta }\left(\frac{1}{E[R_2]}\right),$$
the average delay can be bounded as
$$\mathrm{\Omega }\left(\left(\genfrac{}{}{0pt}{}{N}{N/2}\right)\frac{G}{\mathrm{log}\mathrm{log}G}\right)=D=O\left(\left(\genfrac{}{}{0pt}{}{N}{N/2}\right)G\right).$$
Using Stirlingโs formula, we obtain bounds on the average delay as stated in (21).
## Appendix H Proof of Theorem 10
The average throughput of the multi-group cooperation scheme is given by
$$R_{tot}=E\left[\underset{g=1}{\overset{G}{\mathrm{max}}}\left\{\left(\frac{N}{2}\right)\mathrm{min}\{R_{s1}^g,R_{s2}^g\}\right\}\right].$$
Since $`E[\mathrm{max}\{X_1,\mathrm{},X_n\}]E[X_1]`$, we have
$$R_{tot}E\left[\left(\frac{N}{2}\right)\mathrm{min}\{R_{s1}^1,R_{s2}^1\}\right].$$
The expression on the right is the average throughput of the single group cooperation scheme described in Section 3.3. Using the results of Theorem 8, we have
$$R_{tot}=\mathrm{\Omega }(N).$$
The average throughput can be upper bounded by using the fact that
$$E\left[\underset{g=1}{\overset{G}{\mathrm{max}}}\left\{\left(\frac{N}{2}\right)\mathrm{min}\{R_{s1}^g,R_{s2}^g\}\right\}\right]E\left[\underset{g=1}{\overset{G}{\mathrm{max}}}\frac{NR_{s1}^g}{2}\right]=\left(\frac{N}{2}\right)E\left[\mathrm{log}\left(1+\underset{g=1}{\overset{G}{\mathrm{max}}}|h_{\pi (\frac{N}{2}+1)}^g|^2P\right)\right].$$
The expression on the right is the average throughput of the best among median users scheme proposed earlier. Using the results of Theorem 9, we get
$$R_{tot}=O(N\mathrm{log}\mathrm{log}G).$$
$$\mathrm{\Omega }(N)=R_{tot}=O\left(N\mathrm{log}\mathrm{log}G\right).$$
We now determine the average delay of the multi-group cooperation scheme. To implement this scheme, the base station needs to maintain $`G`$ queues, one for each group. At the beginning of each time slot, the base station selects a group according to condition (22). Since we consider a symmetric scenario, the probability that the base station chooses any particular group is $`(1/G)`$. The information transmitted by the base station in the first half of each time slot reaches all the $`N`$ users in the selected group at the end of that time slot. Hence the average delay for transmitting a packet in the multi-group cooperation scheme is given by
$$D=G(T_c+\mu S)=\mathrm{\Theta }\left(\frac{G}{E[R]}\right),$$
where $`R_{tot}=(NE[R])/2`$. Hence the average delay of the multi-group cooperation scheme can be bounded as
$$\mathrm{\Omega }\left(\frac{G}{\mathrm{log}\mathrm{log}G}\right)=D=O(G).$$
## Appendix I Proof of Lemma 11
From the results on extreme order statistics in , we know that
$$\frac{|\chi _{min}|^2}{b_N}W\text{in distribution},$$
where $`W`$ has a Weibull type distribution and $`b_N`$ satisfies $`F(b_N)=\frac{1}{N}`$, which implies
$$1e^{Lb_N}\left(\underset{k=0}{\overset{L1}{}}\frac{(Lb_N)^k}{k!}\right)=\frac{1}{N}.$$
Using Taylorโs theorem, we get for some $`0<\gamma _N<1`$
$$1e^{Lb_N}(e^{Lb_N}\frac{e^{\gamma _NLb_N}(Lb_N)^L}{L!})=\frac{1}{N}\frac{e^{(1\gamma _N)Lb_N}(Lb_N)^L}{L!}=\frac{1}{N}.$$
Taking $`\mathrm{log}(.)`$ on both sides, we get
$$(1\gamma _N)Lb_NL\mathrm{log}b_N=\mathrm{log}N+L\mathrm{log}L\mathrm{log}(L!).$$
Since $`|\chi _{min}|^2|\chi _1|^2=\mathrm{\Theta }(1)`$, we know that $`b_N=O(1)`$ and hence the $`\mathrm{log}b_N`$ term dominates the left hand side of the above expression. Thus we have
$$b_N=\mathrm{\Theta }\left(N^{\frac{1}{L}}\right).$$
$$N^{\frac{1}{L}}|\chi _{min}|^2kW\text{in distribution, for some constant }k>0\text{.}$$
Since $`E\left[|\chi _{min}|^2\right]E\left[|\chi _1|^2\right]<\mathrm{}`$, we can use the result in Theorem 2.1 of to conclude that
$$N^{\frac{1}{L}}E\left[|\chi _{min}|^2\right]kE[W]=\mathrm{\Theta }(1).$$
$$E\left[|\chi _{min}|^2\right]=\mathrm{\Theta }\left(N^{\frac{1}{L}}\right).$$
The average throughput of the worst user scheme can now be upper bounded using Jensenโs inequality as follows
$$R_{tot}=NE\left[\mathrm{log}\left(1+|\chi _{min}|^2P\right)\right]N\mathrm{log}\left(1+E\left[|\chi _{min}|^2\right]P\right).$$
$$R_{tot}=O\left(N^{\left(\frac{L1}{L}\right)}\right).$$
(57)
We lower bound the average throughput of the worst user scheme as follows
$$R_{tot}=N_0^{\mathrm{}}\mathrm{log}(1+xP)dF_{min}(x)N_{b_N}^{\mathrm{}}\mathrm{log}(1+xP)dF_{min}(x).$$
$$R_{tot}N\mathrm{log}\left(1+b_NP\right)\left[1F_{min}(b_N)\right],$$
where
$$F_{min}(x)=1(1F(x))^N.$$
Using the fact that $`F(b_N)=\frac{1}{N}`$, we get
$$F_{min}(b_N)=1(1F(b_N))^N=1\left(1\frac{1}{N}\right)^N=1e^{N\mathrm{log}\left(1\frac{1}{N}\right)}.$$
$$F_{min}(b_N)=1e^1\left(1+O\left(\frac{1}{N}\right)\right).$$
$$R_{tot}N\mathrm{log}\left(1+b_NP\right)\left[e^1+O\left(\frac{1}{N}\right)\right]=\mathrm{\Theta }\left(N\mathrm{log}\left(1+N^{\frac{1}{L}}P\right)\right).$$
$$R_{tot}=\mathrm{\Omega }\left(N^{\left(\frac{L1}{L}\right)}\right).$$
Combining this with the upper bound in (57), we get
$$R_{tot}=\mathrm{\Theta }\left(N^{\left(\frac{L1}{L}\right)}\right).$$
## Appendix J Proof of Lemma 12
From the results on extreme order statistics in , we know that
$$\frac{|\chi _{max}|^2a_N}{b_N}W\text{in distribution},$$
where $`W`$ has a Gumbel distribution and $`a_N`$ and $`b_N`$ satisfy
$$F(a_N)=1\frac{1}{N}\text{and}b_N=\frac{1}{Nf(a_N)},$$
where $`f(.)`$ denotes the probability density function obtained from (25). Now
$$F(a_N)=1\frac{1}{N}\frac{e^{La_N}(La_N)^{(L1)}}{(L1)!}(1+O\left(\frac{1}{a_N}\right))=\frac{1}{N}.$$
Taking $`\mathrm{log}(.)`$ on both sides and simplifying, we get
$$La_N(L1)\mathrm{log}a_N=\mathrm{log}N+(L1)\frac{1}{2}\mathrm{log}(L1)+K.$$
$$a_N=\frac{\mathrm{log}N+(L1)\mathrm{log}\mathrm{log}N}{L}+O(\mathrm{log}\mathrm{log}N).$$
Since
$$f(a_N)=\frac{Le^{La_N}(La_N)^{(L1)}}{(L1)!}=\mathrm{\Theta }\left(\frac{1}{N}\right),$$
we have $`b_N=C=\mathrm{\Theta }(1)`$. Thus
$$|\chi _{max}|^2\left(\frac{\mathrm{log}N+(L1)\mathrm{log}\mathrm{log}N}{L}+O(\mathrm{log}\mathrm{log}N)\right)CW\text{in distribution}.$$
Using Chebyshevโs inequality, it is easy to show that
$$\frac{|\chi _{max}|^2}{\left(\frac{\mathrm{log}N+(L1)\mathrm{log}\mathrm{log}N}{L}\right)}1\text{in probability}.$$
Since any Chi-squared random variable with $`2L`$ degrees of freedom can be expressed as the sum of $`L`$ exponential i.i.d random variables, we have
$$|\chi _{max}|^2=\underset{i=1}{\overset{N}{\mathrm{max}}}\left\{\frac{Z_1^i+\mathrm{}+Z_L^i}{L}\right\}\underset{i=1}{\overset{N}{\mathrm{max}}}Z_1^i,$$
where $`Z_j^i`$โs are exponential random variables with unit mean. Hence
$$E\left[\frac{|\chi _{max}|^2}{\left(\frac{\mathrm{log}N+(L1)\mathrm{log}\mathrm{log}N}{L}\right)}\right]\frac{E\left[\mathrm{max}_{i=1}^NZ_1^i\right]}{\left(\frac{\mathrm{log}N+(L1)\mathrm{log}\mathrm{log}N}{L}\right)}\frac{k\mathrm{log}N}{\left(\frac{\mathrm{log}N+(L1)\mathrm{log}\mathrm{log}N}{L}\right)}kL<\mathrm{}.$$
Thus we can apply the Dominated Convergence Theorem to get
$$E\left[\frac{|\chi _{max}|^2}{\left(\frac{\mathrm{log}N+(L1)\mathrm{log}\mathrm{log}N}{L}\right)}\right]1E[|\chi _{max}|^2]=\mathrm{\Theta }\left(\frac{\mathrm{log}N+(L1)\mathrm{log}\mathrm{log}N}{L}\right).$$
Using Jensenโs inequality, we get
$$R_{tot}=E\left[\mathrm{log}\left(1+|\chi _{max}|^2P\right)\right]\mathrm{log}\left(1+E\left[|\chi _{max}|^2\right]P\right).$$
$$R_{tot}=O\left(\mathrm{log}\left(1+\frac{\mathrm{log}N+(L1)\mathrm{log}\mathrm{log}N}{L}\right)\right).$$
(58)
We can lower bound the average throughput of the best user scheme as follows
$$R_{tot}=_0^{\mathrm{}}\mathrm{log}(1+xP)dF_{max}(x)_{a_N}^{\mathrm{}}\mathrm{log}(1+xP)dF_{max}(x).$$
$$R_{tot}\mathrm{log}\left(1+a_NP\right)\left[1F_{max}(a_N)\right],$$
where
$$F_{max}(x)=(F(x))^N.$$
Using the fact that $`F(a_N)=1\frac{1}{N}`$, we get
$$F_{max}(a_N)=(F(a_N))^N=\left(1\frac{1}{N}\right)^N=e^1\left(1+O\left(\frac{1}{N}\right)\right).$$
$$R_{tot}\mathrm{log}\left(1+a_NP\right)\left[1e^1+O\left(\frac{1}{N}\right)\right]=\mathrm{\Theta }\left(\mathrm{log}\left(1+a_NP\right)\right).$$
$$R_{tot}=\mathrm{\Omega }\left(\mathrm{log}\left(1+\frac{\mathrm{log}N+(L1)\mathrm{log}\mathrm{log}N}{L}\right)\right).$$
Combining this with the upper bound in (58), we get
$$R_{tot}=\mathrm{\Theta }\left(\mathrm{log}\left(1+\frac{\mathrm{log}N+(L1)\mathrm{log}\mathrm{log}N}{L}\right)\right).$$
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# Dynamical and Pressure Structures in Winds with Multiple Embedded Evaporating Clumps I. 2D Numerical Simulations
## 1 Introduction
Starbursts occur in regions where clumps of cool, dense, molecular material are embedded in a far hotter and more diffuse external medium. An understanding of the responses of such regions to winds impacting on them is central to the investigation of feedback in star and galaxy formation. The same type of structure exists in, for example, planetary nebulae, supernova remnants, Hii regions, and the interstellar medium of the Galactic Center. Clumps (also referred to as clouds or globules) may either accumulate material from the surrounding medium, and thus increase in mass, or may lose material to the external medium, and eventually be destroyed. In the latter case, mass loss can occur through hydrodynamic ablation, or thermal or photoionized evaporation. The diffuse medium is often in motion relative to these clouds, and the nature of this interaction is of wide importance. For instance, in the context of a starburst superwind, the observed X-ray emission will depend on the character of this interaction.
While there have been many studies of the interaction of a dense, cold cloud with a tenuous flow (e.g., Klein, McKee & Colella, 1994; Mac Low et al., 1994; Xu & Stone, 1995; Gregori et al., 2000; Lim, Hartquist & Williams, 2001), there have been few calculations of the interaction between multiple clouds and a flow (e.g., Jun, Jones & Norman, 1996; Poludnenko, Frank & Blackman, 2002). While the latter have recently been supplemented by some wonderful laser experiments (Poludnenko et al. (2004); see also Klein et al. (2003)), our understanding of such interactions is still developing.
A limitation of previous models is that the clouds have been modelled as single phase entities, and the density contrast between the cloud(s) and the flow has typically been taken to be of the order of $`10^2`$ (for numerical reasons). The simulated clouds then have such short lifetimes that they are unable to significantly โmassloadโ the flow. In reality, the astrophysical clouds which are of interest are cold and molecular, and have much larger density contrasts. In this case, the time required to destroy the clouds is much longer than other relevant time-scales and the rate of mass-loss from the cloud can be assumed constant. In this limit the mass-loss may significantly massload the flow, with the injected material occupying a cone with a sizeable opening angle if the wind is hypersonic, or being confined to a long, thin tail when the wind is transonic (Dyson, Hartquist & Biro, 1993; Falle et al., 2002). A short review of tail formation by this, and other processes, is given in Dyson (2003).
The nature of the interaction of a flow with a large group of clouds can differ substantially from that which occurs with a single cloud. If the diffuse medium surrounding embedded molecular clouds is flowing supersonically, then the clouds are likely to be destroyed by hydrodynamical ablation. Mass injection into the flow due to the destruction of the clouds may sometimes greatly enhance the thermal pressure of a flow at the expense of the flowโs ram pressure. The properties of the flow may then be much more conducive to the survival of clouds further downstream, and if the flow is slowed and pressurized enough, it may even induce their collapse and hence trigger new star formation. This process could be a central mechanism for feedback in the interstellar medium (e.g., in starburst regions).
The main features which we might expect to see in the interaction between multiple clouds and a tenuous supersonic flow are shown schematically in Fig. 1. Individual bowshocks form around those clouds furthest upstream and merge at some point downstream. While much of the material in the flow is likely to remain supersonic in this region, further encounters with clouds and the creation of additional bowshocks means that the flow will gradually slow, pressurize, and become subsonic. New star formation may occur in this part of the flow. Other clouds may continue to lose mass, but since they interact with a subsonic flow their injected mass will be confined to long thin tails. Such tails have only a small cross-section to the oncoming flow, and do not greatly impede it. The flow may therefore accelerate between the clouds as this region effectively becomes more porous, and may become supersonic again. If more clouds are encountered further downstream the whole scenario may repeat. One issue is whether such a system would reach a steady state, or whether it would flicker due to mass pickup by the diffuse flows being less effective when the flow is transonic compared to when it is supersonic.
In this paper we use hydrodynamical calculations to investigate the nature of the interaction between a tenuous flow and a number of mass sources. In Section 2 we describe the problem and the basic assumptions used. The numerical results for single to multiple mass sources interacting with both a hypersonic and transonic wind are presented in Section 3. Section 4 contains our conclusions and ideas for future work.
## 2 The Model
As our investigation is still in its early stages, in this study we concern ourselves with only the general properties of the interaction of a wind with injected material. We do not, therefore, attempt to model the detail of how material is injected. Instead, we follow the approach in Falle et al. (2002) by assuming a uniform rate of mass injection within a given radius. This is simple to implement, but has the disadvantage that the flow from the injection region is isotropic. The mass loss from a cloud will not be isotropic if it is caused by hydrodynamic ablation by an incident wind, or by photoevaporation by radiation from a nearby star, but we show in Section 3.1.1 that the effect of asymmetrical mass loss is small at large distances from the cloud. Therefore, the actual details of the mass injection process are relatively unimportant. We also emphasize that the boundary of this injection region is not meant to be the boundary of a cloud. In fact, the cloud could be much smaller.
In this paper we are concerned with the interaction of a flow incident on several mass sources. Three dimensional calculations are required if the sources are spherical. To reduce the computational cost we restrict ourselves to two dimensional simulations where the sources are cylindrical. While we expect some differences between calculations performed in 2D and 3D, at this stage we can still gain important insight from less computationally demanding 2D simulations.
We investigate the simplest case in which the incident wind behaves adiabatically and the injected gas remains isothermal. This is motivated by the fact that if mass injection is due to photoevaporation, the injected temperature will be $`10^4`$ K, which is comparable to that of a wind whose temperature is determined by photoionization - the wind, however, can be shocked to much higher temperatures, which is the reason why we can see tails from mass sources in many astrophysical media<sup>1</sup><sup>1</sup>1If the Mach number of the wind were lower, the density contrast between the injected material and the wind would be reduced, in which case the described dichotomy is a poorer approximation.. To ensure the above behaviour, we use an advected scalar, $`\alpha `$, which is unity in the injected gas and zero in the ambient gas. The source term in the energy equation is then
$$K\alpha \rho (T_0T),$$
(1)
where $`\rho `$ and $`T`$ are the local mass density and temperature, and where $`K`$ is large enough that the temperature always remains close to the equilibrium temperature, $`T_0`$, in the injected gas. Inside the injection region we add an extra energy source so that the gas is injected with temperature unity (see Falle et al. (2002) for further details).
The calculations reported in this paper use Cobra, a 2nd order accurate code with adaptive mesh refinement (AMR). Cobra uses a hierarchy of grids $`G^0\mathrm{}G^N`$ such that the mesh spacing on grid $`G^n`$ is $`\mathrm{\Delta }x_0/2^n`$. Grids $`G^0`$ and $`G^1`$ cover the whole domain, but the finer grids only exist where they are needed. The solution at each position is calculated on all grids that exist there, and the difference between these solutions is used to control refinement. In order to ensure Courant number matching at the boundaries between coarse and fine grids, the time step on grid $`G^n`$ is $`\mathrm{\Delta }t_0/2^n`$ where $`\mathrm{\Delta }t_0`$ is the time step on $`G^0`$. Such a hierarchical grid structure not only improves the efficiency by confining the fine grids to where they are needed, but it also makes it possible to use a full approximation multigrid algorithm to accelerate the convergence to the steady state (see, e.g., Brandt, 1977). Further details of Cobra can be found in Falle & Giddings (1993).
## 3 Results
### 3.1 Single and Twin Mass Sources
Our first investigations are of one and two mass sources interacting with an ambient flow. So that the ambient wind does not affect the flow in the injection region, we must ensure that the ram pressure of the injected material at the boundary of the injection region is larger than that in the wind, i.e. we require
$$a^2\rho _\mathrm{s}=\frac{ar_\mathrm{c}Q}{3}\rho _\mathrm{w}v_\mathrm{w}^2,$$
(2)
where $`r_\mathrm{c}`$ is the radius of the injection region, $`a`$ is the flow speed of injected material at $`r=r_\mathrm{c}`$ and $`\rho _\mathrm{s}`$ its density at this point, $`Q`$ is the mass injection rate per unit volume within $`r=r_\mathrm{c}`$, and $`\rho _\mathrm{w}`$ and $`v_\mathrm{w}`$ are the wind density and velocity (cf. Equation 5 in Falle et al. (2002)). We again emphasize that the cloud radius could actually be much smaller than $`r_\mathrm{c}`$.
#### 3.1.1 Hypersonic wind
We choose units such that $`r_\mathrm{c}=1`$, $`a=1`$, $`\rho _\mathrm{w}=1`$, and set $`v_\mathrm{w}=20`$ and $`T_0=1`$, which correspond to an external isothermal Mach number of 20. Equation 2 then implies $`Q800`$ in order to prevent the wind from penetrating the mass injection region, although Falle et al. (2002) note that in the supersonic case a better estimate is obtained by replacing $`\rho _\mathrm{w}v_\mathrm{w}^2`$ by the pressure behind a stationary normal shock in the wind. This gives
$$Q\frac{4}{\gamma +1}\frac{\rho _\mathrm{w}v_\mathrm{w}^2}{ar_\mathrm{g}}=600$$
(3)
for $`\gamma =5/3`$. We set $`Q=700`$ in order to ensure that the interaction occurs slightly outside the mass-injection region. As in Falle et al. (2002), we set the coefficient $`K`$ to 50, which is large enough to ensure that the temperature in the injected gas stays close to unity.
The computational domain is $`0x50`$, $`50y20`$ with the mass injection region initially centered at the origin. 5 grid levels $`G^0\mathrm{}G^4`$ were used, with $`G^0`$ being $`50\times 70`$ and $`G^4`$ $`800\times 1120`$. The diameter of the mass injection region is 32 cells on the $`G^4`$ grid. Artificial dissipation is added to the simulations in order to eliminate the โcarbuncle effectโ<sup>2</sup><sup>2</sup>2This is an unphysical distortion of a shock front that is partially aligned with the grid (Quik, 1994). It can be cured by adding a small amount of artificial viscosity to the Riemann solver (e.g., Falle, Komissarov & Joarder, 1998). This artificial dissipation is quite separate from the turbulence model noted in Section 3.1.2, and its influence is confined to the grid scale. and to damp Kelvin-Helmholtz instabilities due to velocity shear at the contact discontinuity. This allows a steady-state solution.
Fig. 2 shows the density, pressure, and y-velocity, together with the regions occupied by the finest grid. The wind is decelerated by a bow shock, which stands off from the center of the mass source by a distance of approximately 8 units. A contact discontinuity separates the wind from the injected material. As noted in Falle et al. (2002) for the case of a spherical mass source, the injected material is not confined to a long thin tail. Instead the contact discontinuity has a half-width of $`20`$ at $`y=50`$, which is much wider than the mass source. It is also much wider than the equivalent width observed at the same distance downstream from a spherical mass source, due to the reduced divergence in the 2D simulation. It appears that the contact discontinuity does not reach an asymptotic off-axis distance far downstream, but rather tends towards a finite opening angle (this can be further discerned in Fig. 12). Though not shown here, the opening angle is dependent on the Mach number of the wind - the injected material is more confined at lower Mach numbers, though the opening angle of the bowshock is greater. A reverse shock surrounds the mass source, and delineates the position at which the isotropic injected material feels the presence of the ambient wind.
In Fig. 3 we show density plots from simulations where the density, velocity, and pressure of the injected material is specified around the edge of the injection region. This approach allows us to investigate the effect of non-isotropic mass injection on the interaction with the ambient wind. In the simulation shown in the left panel of Fig. 3 we set $`\rho _\mathrm{s}=350`$ and $`a=1`$, and keep the other parameters as before. The total mass-injection rate is the same as that in Fig. 2, though the energy injection is slightly different. The latter means that there are slight differences between the density plots in Figs. 2 and 3. In the middle panel of Fig. 3 we keep the same overall mass injection rate, but vary the density at the edge of the injection region according to the prescription
$$\rho _\mathrm{s}=\rho _0(1\mathrm{\Omega }^2\mathrm{sin}^2\theta ),$$
(4)
where $`\rho _0`$ is a normalization factor, $`\theta `$ is the angle between the radial vector from the injection region and the ambient flow ($`\theta =\pi `$ on the upstream surface, and $`\theta =0`$ on the downstream surface), and $`\mathrm{\Omega }`$ sets the degree of anisotropy. We set $`\mathrm{\Omega }=0.9`$, so that the mass injection rate at the upstream and downstream surfaces is $`5\times `$ that at $`\theta =\pi /2`$. While there are differences in the morphology of the interaction close to the injection region (e.g., the bow shock stands further off, and the shape of the reverse shock reflects the latitudinal variation of the mass and energy injection), the large scale features are remarkably unchanged. In the right panel of Fig. 3 the injection rate is highest at the upstream surface and declines smoothly towards the downstream surface, according to the prescription
$$\rho _\mathrm{s}=\rho _0[1+\mathrm{\Omega }\mathrm{sin}(\theta +\pi /2)].$$
(5)
This prescription may be expected to be more representative of reality. With $`\mathrm{\Omega }=0.7`$, the mass injection rate at the upstream surface is $`5\times `$ that at the downstream surface. Again we see that the large scale features are very similar, though as expected the bowshock stands slightly further off than in the other two cases.
In Fig. 4 we show the density from calculations where the mass source is moved off-axis, which simulates the interaction of a wind with 2 identical sources. The separation between the two sources increases from 12, to 24, to 48, to 96 units. In the top left panel of Fig. 4 the two sources are surrounded by a global contact discontinuity, and not individual ones. The bow shock is located further upstream compared to that in Fig. 2, and is also global in the sense that it envelops the two mass sources, as opposed to individual shocks surrounding each source. When the mass sources are still fairly close together, the flow between them is at a higher pressure than the corresponding flow around the outside edge of the interaction. This causes the flow around each mass source to angle away from each other, and is easily identified by the tilt of the reverse shock around each source relative to the flow of the ambient wind.
As the mass sources are separated, first the global contact discontinuity splits into individual contacts around each source, and then the global bow shock separates into individual bow shocks. The reverse shock around each mass source is once again aligned with the oncoming wind, and the tail of injected material is initially symmetrical and unaffected by the presence of the other mass source. However, some interaction occurs further downstream where the bowshocks interact. At this point a reflected shock is formed, and further downstream this deflects the contact discontinuity and the tail of injected material away from the axis of symmetry. The simulations shown in Fig. 4 used 4 grid levels $`G^0\mathrm{}G^3`$. The diameter of the mass injection region is 16 cells on the $`G^3`$ grid. For the models with separations of 12, 24, and 48 units, the $`G^3`$ grid is $`400\times 560`$, and the computational domain is $`0x50`$, $`50y20`$. The model with a separation of 96 units has a $`G^3`$ grid of $`768\times 960`$, and a computational domain which encompasses $`0x96`$, $`100y20`$.
#### 3.1.2 Transonic wind
In many situations it is usual for a wind to encounter a large number of mass sources. In such circumstances the Mach number of the wind is driven towards unity (Hartquist et al., 1986), and it is therefore reasonable to look at the interaction of mass sources with a transonic wind. In this case we set $`r_\mathrm{c}=1`$, $`a=1`$, and $`T_0=1`$ as before but now $`\rho _\mathrm{w}=10^3`$, $`v_\mathrm{w}=40.825`$, giving a Mach number of unity in the undisturbed wind. Equation 2 gives $`Q10/3`$, which is a good estimate for this case since there is no shock upstream of the injection region. We therefore set $`Q=10/3`$, and the coefficient $`K`$ to $`10^2`$ (cf. Falle et al., 2002).
Simulations with a transonic wind pose additional difficulties compared to the hypersonic wind case. First, in order for the boundaries to have no effect on the solution, the computational domain has to be very large, particularly since our 2D simulations have reduced divergence compared to the case of spherical mass injection regions. To satisfy this condition we find that we need $`0x320`$, $`416y256`$, which would have made the calculation extremely expensive if an AMR code were not employed. Second, the velocity shear between the injected material and the wind is so extreme in the transonic case that the wind flow separates and produces a turbulent wake downstream of the interaction region. Calculations based on the Euler equations cannot adequately describe such turbulence, and the only viable option is to use a turbulence model. While there are many possibilities, we use a simple $`kฯต`$ model as used in Falle et al. (2002).
The purpose of the subgrid turbulence model is to emulate a high Reynolds number flow. It does so by including equations for the turbulent energy density and dissipation rate and using these to calculate viscous and diffusive terms. The resulting solution should be an approximation to the mean flow. The turbulent mixing of the injected gas with the original flow due to shear instabilities is modelled by the diffusive terms in the equations. Since the turbulent viscosity computed from the subgrid model depends upon the local solution, it is not very meaningful to talk about an effective Reynolds number. For example, the turbulent viscosity is largest in shear layers and essentially vanishes in regions with little shear. The model has been calibrated by comparing the computed growth of shear layers with experiments (Dash & Wolf, 1983). The model also assumes that the real Reynolds number is very large, which is the case in astrophysical flows, and that the turbulence is fully developed. Although not entirely satisfactory, such a model gives a much more realistic result than simply using grid viscosity since in that case the size of the shear instabilities are determined by the numerical resolution. Further details can be found in Falle (1994). The effect of the turbulence model is illustrated in Fig. 5.
We use 9 levels of grid refinement for our transonic simulations, $`G^0\mathrm{}G^8`$, with $`G^0`$ being $`10\times 21`$ and $`G^8`$ $`2560\times 5376`$. The diameter of the mass injection region is 16 cells on the $`G^8`$ grid. The interaction is very different from the hypersonic case, as can be seen from Fig. 6. Instead of a bow shock there is a bow wave upstream of the mass source, whose amplitude falls off as $`1/r`$. A very weak tail shock occurs in the wind downstream of the mass source, and this is aligned much more parallel to the ambient flow than when the mass source is spherical (cf. Falle et al., 2002). The injected material remains in rough pressure equilibrium with the wind, and eventually becomes confined to a tail whose width is of about the same order as the injection region.
In Fig. 7 we show the density, pressure, and y-velocity from a calculation where the mass source is moved off-axis, as we again simulate the interaction between a wind and 2 identical sources. The separation between the two sources is 48 units. We immediately see that the tail produced behind each mass source is strongly curved towards the other, producing a narrow channel between the tails. However, it is clear that the tails are initially tilted away from each other, as shown by the position of the reverse shock around the mass source. The widening of the channel between the two mass sources, which is enhanced by the curvature of the inner side of the contact discontinuity, causes the pressure within the channel to drop via the Bernoulli effect as the wind accelerates and becomes supersonic. The fall in pressure relative to that on the outside edge of the tail results in the tail curving towards the axis of symmetry, and the channel is then closed off. A shock in the channel just above this point slows the accelerated wind. There is no sign of any tail shock in the downstream wind on the outer side of the tail.
Further simulations reveal that the curvature of the tails towards each other decrease when the separation between the mass sources is increased. At large enough separations there is no interaction between the mass sources and the tails remain perfectly straight.
### 3.2 Multiple Sources in a Hypersonic Wind
In this section we investigate the interaction between a hypersonic flow and a group of 5 mass sources (10 with the imposed symmetry), and use the turbulence model noted in Section 3.1.2 for these calculations.
In the first simulation the sources are randomly distributed within a circular region of diameter 160 units. This diameter is increased in subsequent simulations, in order to explore differences between an interaction which is dominated by the mass injection to one which is dominated by the wind, though the relative positions of the sources remain the same. We define the parameter $`\chi =\dot{M}_\mathrm{c}/\dot{M}_\mathrm{w}`$, where $`\dot{M}_\mathrm{c}`$ is the combined mass injection rate of the sources and $`\dot{M}_\mathrm{w}`$ is the mass flux in the wind through a suitably chosen region. When we change the size of the region in which the mass sources are distributed, we alter the value of $`\chi `$ since $`\rho _\mathrm{w}`$ and $`v_\mathrm{w}`$ are kept constant. The mass injection regions have diameters of 8 cells on the finest grid in each of the models presented in this section.
The positions of the 5 sources in our first simulation are $`(x,y)=(33.72,57.12)`$, (48.69,-13.70), (39.07,2.54), (24.68,55.46), and (7.89,-14.51). The best estimate for $`\dot{M}_\mathrm{w}`$ is obtained from the evaluation of the flow rate through a region bounded by the symmetry axis and the source region whose $`x`$-coordinate is furthest from this. $`D`$ is set to the value of this $`x`$-coordinate. Thus $`\dot{M}_\mathrm{w}=\rho _\mathrm{w}v_\mathrm{w}D=20\times 1\times 48.69=973.8`$. Because we are adding mass and energy to the wind, we use a slightly higher value of $`Q`$ in these simulations to ensure that the flow does not enter any of our source regions ($`Q=880`$). Therefore, $`\chi =5\pi \times 880/973.814`$ and we expect the interaction to be injection-dominated. We use 5 grid levels, $`G^0\mathrm{}G^4`$ with $`G^0`$ being $`40\times 60`$ and $`G^4`$ $`640\times 960`$. The computational domain is $`0x160`$, $`160y80`$.
In Fig. 8 we show the density, pressure, y-velocity, and the advected scalar of this model. A global bow shock exists around the group of mass sources and the region between the sources is filled with high pressure, low Mach number gas. The shape of the global bow shock is to some extent determined by the positions of the individual sources, and has an opening angle which is much wider than that obtained when there are only one or two sources (cf. Section 3.1.1). Downstream of the mass sources the injected material is accelerated by a pressure gradient in a manner similar to that of a superwind. The shape and position of the reverse shock around each mass source is dependent on the local flow conditions, and in some cases is almost spherical. Close up images of some of the injection regions are shown in Fig. 9. The bottom right panel of Fig. 8 reveals that the group of mass sources is largely impervious to the oncoming wind. The mass source which is furthest upstream is somewhat akin to a continental divide in the sense that it splits the wind flow to the left or right. Since we impose symmetry at $`x=0`$ a high pressure region of shocked ambient wind is formed - wind in this region is able to percolate through the region of mass sources, and is the only part of the oncoming flow which is able to do so.
If the distances between the mass sources are increased by a factor of 4, $`\chi `$ is reduced to $`3.5`$. The interaction between the wind and the injected material is still dominated by the mass sources, as shown in Figs. 10 and 11, but to a lesser extent than the simulation shown in Fig. 8. The mass source furthest upstream again acts like a continental divide. However, the flow in the vicinity of this source is identical to the single source case, being unaltered by the complexities of the interaction further downstream. Specifically, the reverse shock around the mass source is not tilted or compressed. Where the interaction with $`\chi 14`$ is characterized by a large region of subsonic flow, the interaction with $`\chi 3.5`$ has a complex structure with multiple shocks, as is readily apparent in Fig. 10. This is due to the fact that the distances between the mass sources are such that pressure gradients in the vicinity of each source accelerate the flow to supersonic velocities before additional mass sources are encountered further downstream. In addition to a region near the axis of symmetry, the advected scalar reveals that the wind is able to force its way between the mass sources in two distinct streams, which become diluted by injected material along their length. In this simulation we used 7 grid levels, $`G^0\mathrm{}G^6`$ with $`G^0`$ being $`40\times 60`$ and $`G^6`$ $`2560\times 3840`$, spanning a computational domain of $`0x640`$, $`640y320`$.
An increase in the separation of the mass sources by another factor of 4 reduces $`\chi `$ to approximately unity. We now expect the wind to begin to force its way through the region of mass sources, and this is demonstrated in Fig. 12, a simulation making use of 9 grid levels, $`G^0\mathrm{}G^8`$ with $`G^0`$ being $`40\times 60`$ and $`G^8`$ $`10240\times 15360`$ (157 million cells equivalent). The computational domain spans the region $`0x2560`$, $`2560y1280`$. The distances between each mass source are now so great that there is very little interaction between them, the main difference being that the sources which are furthest downstream are interacting with a supersonic flow whose properties have been somewhat modified by the action of the sources further upstream.
The density and Mach number averaged across $`x`$ in the disturbed flow as a function of $`y`$ for each of the three simulations discussed in this section are shown in Fig. 13. The Mach number is mass averaged and not volume averaged. We see that when the rate of mass injection dominates the mass flux of the wind (i.e. $`\chi 1`$) the average density across the width of the interaction region is very high, being of order 200 times the ambient density of the wind over a large volume of the region that contains the mass sources. The corresponding Mach number in this case is typically 0.5. The density decreases and the Mach number increases downstream of the region containing the mass sources, and the flow passes through a sonic point at roughly the same $`y`$-coordinate as that of the most downstream mass source. The density and Mach number profiles are fairly smooth, though there are some features with small spatial scales.
As $`\chi `$ decreases, the mass sources have a much more localized effect on the flow. When $`\chi =3.5`$, we see that the flow between the most upstream mass source (at $`y222`$) and its nearest companion (at $`y10`$) becomes supersonic, and thus knows nothing of the presence of this second source prior to the bowshock around it<sup>3</sup><sup>3</sup>3This is in marked contrast to the simulation with $`\chi =14`$ where a high pressure region moves upstream and compresses the reverse shock around the mass source furthest upstream.. The density in this part of the flow thus significantly decreases from its peak post-shock value around the first mass source ($`\rho 200\rho _{\mathrm{amb}}`$) until its encounter with the second mass source ($`\rho 8\rho _{\mathrm{amb}}`$ just ahead of the bowshock). There are a few places in the interaction region where the flow averaged over the width of the interaction is subsonic, and these are associated with local maxima in the corresponding density profiles. Once again, the overall flow passes through a sonic point close to the $`y`$-coordinate of the most downstream mass source. Compared to the simulation with $`\chi =14`$, the density enhancement of the mass-loaded flow is much reduced, and the averaged flow is able to maintain a higher Mach number. The trends which we have noted above continue as $`\chi `$ is made yet smaller. When $`\chi =0.88`$ the individual mass sources appear to act only as localized perturbations on the flow.
The pressure at a given radius from the centre of specific mass sources is shown in Fig. 14. In this figure we have chosen mass sources with the most uniform pressure surroundings in the two models considered. We see that the pressure around a specific mass source tends to be higher when $`\chi `$ is large<sup>4</sup><sup>4</sup>4The exception to this is a small region at an azimuthal angle of $`15^{}`$ around the mass source with position $`(x,y)=(31.57,58.04)`$ in the simulation with $`\chi =3.5`$.. When the reverse shock extends to radii exceeding that at which the pressure profile is obtained, the pressure drops to a low value. In the simulation with $`\chi =3.5`$ we typically find that the pressure is lowest on the downstream side of the mass source, and highest on the upstream side. However, this โmemoryโ of the ambient flow is reduced as $`\chi `$ increases - in the top left panel of Fig. 14 we see that the maximum pressure occurs on the downstream side of the mass source. Another noteable feature is that the pressure as a function of azimuthal angle around the mass source may be fairly constant when $`\chi `$ is large. We anticipate that if the pressure is high and relatively uniform then there will be a greater probability that the clump could collapse and give rise to new star formation, than when the pressure is low, or varies significantly around the clump.
Although we raised the possibility of flickering in Section 1, we see little evidence for this in our models. The horizontal shock in the injected material at $`y10`$ in the $`\chi 3.5`$ model (see Fig. 10 and the bottom left panel of Fig. 11) was observed to display some vertical oscillations, but these may be the result of the system relaxing to a steady state and we have not evolved the simulation long enough to eliminate this possibility. We anticipate that flickering may be more apparent in a simulation where the mass injection rate from each source responds to the local flow conditions.
## 4 Conclusions
The results shown in Section 3 illustrate that a 2D calculation of mass injection from a cylindrical source can produce flow features which are similar to those obtained in an axisymmetric simulation where mass is injected from a spherical source (cf. Falle et al., 2002). For the interaction with a hypersonic wind, these features include the presence of a bow shock and the fact that injected material occupies a downstream region with a width significantly larger than that of the injection region. For the interaction with a transonic wind, the tail produced in the cylindrical case is slightly broader than that produced in the spherical case, but in both cases the injected material occupies a region whose width is considerably less than that obtained when the wind is hypersonic.
When two mass sources are in close proximity, their interaction may affect the shape and alignment of the tail which forms downstream of each source. Deviation in the alignment of the tails is caused by pressure differences either side of each mass source and by interaction with the reflected bow shock. These induce velocity components into the flow perpendicular to the upstream velocity of the ambient wind. A global bow shock envelops the mass sources when they are close to each other, but increasing their separation leads to an individual bow shock around each source.
For the transonic case, the curvature of the tails is also produced by pressure differences. The region between the two tails acts as a narrow channel, and the curvature of the contact discontinuity causes the wind material which travels between the tails to be accelerated. This is accompanied by a pressure drop, and the pressure differences on either side of the tail force the injected material towards the symmetry axis. In both the hypersonic and transonic cases, separating the sources reduces the degree of the interaction effects.
An interesting result is that in the hypersonic case, tails which show deviations from alignment with the upstream wind velocity are pointing away from the mass sources, whereas in the transonic case the tails end up pointing towards each other. While this could be a useful way of determining whether the wind impacting on two mass sources which are close together is hypersonic or transonic, it is possible that in a 3D model the wind between the mass sources would not be accelerated as much, leading to smaller pressure differences on either side of a tail and less significant curvature. The morphology of the tail (broad and โstubbyโ, versus long and thin) is instead a much simpler way to determine the nature of the interaction.
With multiple mass sources in a hypersonic flow, the ability of the wind to punch through the space between the sources depends on the ratio of the mass injection rate from the sources to the mass flux in the wind, which we have called $`\chi `$. When $`\chi `$ is much greater than unity, the sources are an effective barrier, and allow very little of the impacting wind to find a path through them. A global bowshock exists around the sources and the space between the sources is filled with a high pressure, subsonic flow of injected material. When $`\chi `$ is less than or of order unity, the wind is able to force its way inbetween the mass sources, and for the most part this flow is much less pressurized and highly supersonic. We see little evidence for flickering in our current models.
In future work we will consider 3D simulations, different treatments of cooling, and the response of the mass injection rate of each source to the local flow conditions. We will also apply these results to specific objects (e.g., planetary nebulae).
## acknowledgements
JMP would like to thank PPARC for the funding of a PDRA position and current funding from the Royal Society. This research has made use of NASAโs Astrophysics Data System Abstract Service.
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# Seeing stars: Exploiting class relationships for sentiment categorization with respect to rating scales
## 1 Introduction
There has recently been a dramatic surge of interest in sentiment analysis, as more and more people become aware of the scientific challenges posed and the scope of new applications enabled by the processing of subjective language. (The papers collected by Qu, Shanahan, and Wiebe form a representative sample of research in the area.) Most prior work on the specific problem of categorizing expressly opinionated text has focused on the binary distinction of positive vs. negative . But it is often helpful to have more information than this binary distinction provides, especially if one is ranking items by recommendation or comparing several reviewersโ opinions: example applications include collaborative filtering and deciding which conference submissions to accept.
Therefore, in this paper we consider generalizing to finer-grained scales: rather than just determine whether a review is โthumbs upโ or not, we attempt to infer the authorโs implied numerical rating, such as โthree starsโ or โfour starsโ. Note that this differs from identifying opinion strength : rants and raves have the same strength but represent opposite evaluations, and referee forms often allow one to indicate that one is very confident (high strength) that a conference submission is mediocre (middling rating). Also, our task differs from ranking not only because one can be given a single item to classify (as opposed to a set of items to be ordered relative to one another), but because there are settings in which classification is harder than ranking, and vice versa.
One can apply standard $`n`$-ary classifiers or regression to this rating-inference problem; independent work by ?) considers such methods. But an alternative approach that explicitly incorporates information about item similarities together with label similarity information (for instance, โone starโ is closer to โtwo starsโ than to โfour starsโ) is to think of the task as one of metric labeling , where label relations are encoded via a distance metric. This observation yields a meta-algorithm, applicable to both semi-supervised (via graph-theoretic techniques) and supervised settings, that alters a given $`n`$-ary classifierโs output so that similar items tend to be assigned similar labels.
In what follows, we first demonstrate that humans can discern relatively small differences in (hidden) evaluation scores, indicating that rating inference is indeed a meaningful task. We then present three types of algorithms โ one-vs-all, regression, and metric labeling โ that can be distinguished by how explicitly they attempt to leverage similarity between items and between labels. Next, we consider what item similarity measure to apply, proposing one based on the positive-sentence percentage. Incorporating this new measure within the metric-labeling framework is shown to often provide significant improvements over the other algorithms.
We hope that some of the insights derived here might apply to other scales for text classifcation that have been considered, such as clause-level opinion strength ; affect types like disgust ; reading level ; and urgency or criticality .
## 2 Problem validation and formulation
We first ran a small pilot study on human subjects in order to establish a rough idea of what a reasonable classification granularity is: if even people cannot accurately infer labels with respect to a five-star scheme with half stars, say, then we cannot expect a learning algorithm to do so. Indeed, some potential obstacles to accurate rating inference include lack of calibration (e.g., what an understated author intends as high praise may seem lukewarm), author inconsistency at assigning fine-grained ratings, and ratings not entirely supported by the text <sup>1</sup><sup>1</sup>1For example, the critic Dennis Schwartz writes that โsometimes the review itself \[indicates\] the letter grade should have been higher or lower, as the review might fail to take into consideration my overall impression of the film โ which I hope to capture in the gradeโ (http://www.sover.net/~ozus/cinema.htm)..
For data, we first collected Internet movie reviews in English from four authors, removing explicit rating indicators from each documentโs text automatically. Now, while the obvious experiment would be to ask subjects to guess the rating that a review represents, doing so would force us to specify a fixed rating-scale granularity in advance. Instead, we examined peopleโs ability to discern relative differences, because by varying the rating differences represented by the test instances, we can evaluate multiple granularities in a single experiment. Specifically, at intervals over a number of weeks, we authors (a non-native and a native speaker of English) examined pairs of reviews, attemping to determine whether the first review in each pair was (1) more positive than, (2) less positive than, or (3) as positive as the second. The texts in any particular review pair were taken from the same author to factor out the effects of cross-author divergence.
As Table 1 shows, both subjects performed perfectly when the rating separation was at least 3 โnotchesโ in the original scale (we define a notch as a half star in a four- or five-star scheme and 10 points in a 100-point scheme). Interestingly, although human performance drops as rating difference decreases, even at a one-notch separation, both subjects handily outperformed the random-choice baseline of 33%. However, there was large variation in accuracy between subjects. <sup>2</sup><sup>2</sup>2 One contributing factor may be that the subjects viewed disjoint document sets, since we wanted to maximize experimental coverage of the types of document pairs within each difference class. We thus cannot report inter-annotator agreement, but since our goal is to recover a reviewerโs โtrueโ recommendation, reader-author agreement is more relevant. While another factor might be degree of English fluency, in an informal experiment (six subjects viewing the same three pairs), native English speakers made the only two errors.
Because of this variation, we defined two different classification regimes. From the evidence above, a three-class task (categories 0, 1, and 2 โ essentially โnegativeโ, โmiddlingโ, and โpositiveโ, respectively) seems like one that most people would do quite well at (but we should not assume 100% human accuracy: according to our one-notch results, people may misclassify borderline cases like 2.5 stars). Our study also suggests that people could do at least fairly well at distinguishing full stars in a zero- to four-star scheme. However, when we began to construct five-category datasets for each of our four authors (see below), we found that in each case, either the most negative or the most positive class (but not both) contained only about 5% of the documents. To make the classes more balanced, we folded these minority classes into the adjacent class, thus arriving at a four-class problem (categories 0-3, increasing in positivity). Note that the four-class problem seems to offer more possibilities for leveraging class relationship information than the three-class setting, since it involves more class pairs. Also, even the two-category version of the rating-inference problem for movie reviews has proven quite challenging for many automated classification techniques .
We applied the above two labeling schemes to a scale dataset<sup>3</sup><sup>3</sup>3Available at http://www.cs.cornell.edu/People/pabo/movie-review-data as scale dataset v1.0. containing four corpora of movie reviews. All reviews were automatically pre-processed to remove both explicit rating indicators and objective sentences; the motivation for the latter step is that it has previously aided positive vs. negative classification . All of the 1770, 902, 1307, or 1027 documents in a given corpus were written by the same author. This decision facilitates interpretation of the results, since it factors out the effects of different choices of methods for calibrating authorsโ scales. <sup>4</sup><sup>4</sup>4 From the Rotten Tomatoes websiteโs FAQ: โstar systems are not consistent between critics. For critics like Roger Ebert and James Berardinelli, 2.5 stars or lower out of 4 stars is always negative. For other critics, 2.5 stars can either be positive or negative. Even though Eric Lurio uses a 5 star system, his grading is very relaxed. So, 2 stars can be positive.โ Thus, calibration may sometimes require strong familiarity with the authors involved, as anyone who has ever needed to reconcile conflicting referee reports probably knows. We point out that it is possible to gather author-specific information in some practical applications: for instance, systems that use selected authors (e.g., the Rotten Tomatoes movie-review website โ where, we note, not all authors provide explicit ratings) could require that someone submit rating-labeled samples of newly-admitted authorsโ work. Moreover, our results at least partially generalize to mixed-author situations (see Section 5.2).
## 3 Algorithms
Recall that the problem we are considering is multi-category classification in which the labels can be naturally mapped to a metric space (e.g., points on a line); for simplicity, we assume the distance metric $`d(\mathrm{},\mathrm{}^{})=|\mathrm{}\mathrm{}^{}|`$ throughout. In this section, we present three approaches to this problem in order of increasingly explicit use of pairwise similarity information between items and between labels. In order to make comparisons between these methods meaningful, we base all three of them on Support Vector Machines (SVMs) as implemented in Joachimsโ $`\mathrm{SVM}^{light}`$ package.
### 3.1 One-vs-all
The standard SVM formulation applies only to binary classification. One-vs-all (OVA) is a common extension to the $`n`$-ary case. Training consists of building, for each label $`\mathrm{}`$, an SVM binary classifier distinguishing label $`\mathrm{}`$ from โnot-$`\mathrm{}`$โ. We consider the final output to be a label preference function $`\pi ^{\mathrm{ova}}(x,\mathrm{})`$, defined as the signed distance of (test) item $`x`$ to the $`\mathrm{}`$ side of the $`\mathrm{}`$ vs. not-$`\mathrm{}`$ decision plane.
Clearly, OVA makes no explicit use of pairwise label or item relationships. However, it can perform well if each class exhibits sufficiently distinct language; see Section 4 for more discussion.
### 3.2 Regression
Alternatively, we can take a regression perspective by assuming that the labels come from a discretization of a continuous function $`g`$ mapping from the feature space to a metric space. <sup>5</sup><sup>5</sup>5We discuss the ordinal regression variant in Section 6. If we choose $`g`$ from a family of sufficiently โgradualโ functions, then similar items necessarily receive similar labels. In particular, we consider linear, $`\epsilon `$-insensitive SVM regression ; the idea is to find the hyperplane that best fits the training data, but where training points whose labels are within distance $`\epsilon `$ of the hyperplane incur no loss. Then, for (test) instance $`x`$, the label preference function $`\pi ^{\mathrm{reg}}(x,\mathrm{})`$ is the negative of the distance between $`\mathrm{}`$ and the value predicted for $`x`$ by the fitted hyperplane function.
?) used SVM regression to classify clause-level strength of opinion, reporting that it provided lower accuracy than other methods. However, independently of our work, ?) found that applying linear regression to classify documents (in a different corpus than ours) with respect to a three-point rating scale provided greater accuracy than OVA SVMs and other algorithms.
### 3.3 Metric labeling
Regression implicitly encodes the โsimilar items, similar labelsโ heuristic, in that one can restrict consideration to โgradualโ functions. But we can also think of our task as a metric labeling problem , a special case of the maximum a posteriori estimation problem for Markov random fields, to explicitly encode our desideratum. Suppose we have an initial label preference function $`\pi (x,\mathrm{})`$, perhaps computed via one of the two methods described above. Also, let $`d`$ be a distance metric on labels, and let $`nn_k(x)`$ denote the $`k`$ nearest neighbors of item $`x`$ according to some item-similarity function $`\mathrm{๐ ๐๐}`$. Then, it is quite natural to pose our problem as finding a mapping of instances $`x`$ to labels $`\mathrm{}_x`$ (respecting the original labels of the training instances) that minimizes
$$\underset{x\text{test}}{}\left[\pi (x,\mathrm{}_x)+\alpha \underset{ynn_k(x)}{}f(d(\mathrm{}_x,\mathrm{}_y))\mathrm{๐ ๐๐}(x,y)\right],$$
where $`f`$ is monotonically increasing (we chose $`f(d)=d`$ unless otherwise specified ) and $`\alpha `$ is a trade-off and/or scaling parameter. (The inner summation is familiar from work in locally-weighted learning <sup>6</sup><sup>6</sup>6 If we ignore the $`\pi (x,\mathrm{})`$ term, different choices of $`f`$ correspond to different versions of nearest-neighbor learning, e.g., majority-vote, weighted average of labels, or weighted median of labels. .) In a sense, we are using explicit item and label similarity information to increasingly penalize the initial classifier as it assigns more divergent labels to similar items.
In this paper, we only report supervised-learning experiments in which the nearest neighbors for any given test item were drawn from the training set alone. In such a setting, the labeling decisions for different test items are independent, so that solving the requisite optimization problem is simple.
#### Aside: transduction
The above formulation also allows for transductive semi-supervised learning as well, in that we could allow nearest neighbors to come from both the training and test sets. We intend to address this case in future work, since there are important settings in which one has a small number of labeled reviews and a large number of unlabeled reviews, in which case considering similarities between unlabeled texts could prove quite helpful. In full generality, the corresponding multi-label optimization problem is intractable, but for many families of $`f`$ functions (e.g., convex) there exist practical exact or approximation algorithms based on techniques for finding minimum s-t cuts in graphs . Interestingly, previous sentiment analysis research found that a minimum-cut formulation for the binary subjective/objective distinction yielded good results . Of course, there are many other related semi-supervised learning algorithms that we would like to try as well; see ?) for a survey.
## 4 Class struggle: finding a label-correlated item-similarity function
We need to specify an item similarity function $`\mathrm{๐ ๐๐}`$ to use the metric-labeling formulation described in Section 3.3. We could, as is commonly done, employ a term-overlap-based measure such as the cosine between term-frequency-based document vectors (henceforth โTO(cos)โ). However, Table 2 shows that in aggregate, the vocabularies of distant classes overlap to a degree surprisingly similar to that of the vocabularies of nearby classes. Thus, item similarity as measured by TO(cos) may not correlate well with similarity of the itemโs true labels.
We can potentially develop a more useful similarity metric by asking ourselves what, intuitively, accounts for the label relationships that we seek to exploit. A simple hypothesis is that ratings can be determined by the positive-sentence percentage (PSP) of a text, i.e., the number of positive sentences divided by the number of subjective sentences. (Term-based versions of this premise have motivated much sentiment-analysis work for over a decade .) But counterexamples are easy to construct: reviews can contain off-topic opinions, or recount many positive aspects before describing a fatal flaw.
We therefore tested the hypothesis as follows. To avoid the need to hand-label sentences as positive or negative, we first created a sentence polarity dataset <sup>7</sup><sup>7</sup>7Available at http://www.cs.cornell.edu/People/pabo/movie-review-data as sentence polarity dataset v1.0. consisting of 10,662 movie-review โsnippetsโ (a striking extract usually one sentence long) downloaded from www.rottentomatoes.com; each snippet was labeled with its source reviewโs label (positive or negative) as provided by Rotten Tomatoes. Then, we trained a Naive Bayes classifier on this data set and applied it to our scale dataset to identify the positive sentences (recall that objective sentences were already removed).
Figure 1 shows that all four authors tend to exhibit a higher PSP when they write a more positive review, and we expect that most typical reviewers would follow suit. Hence, PSP appears to be a promising basis for computing document similarity for our rating-inference task. In particular, we defined $`\stackrel{}{\mathrm{PSP}(x)}`$ to be the two-dimensional vector $`(\mathrm{PSP}(x),1\mathrm{PSP}(x))`$, and then set the item-similarity function required by the metric-labeling optimization function (Section 3.3) to $`\mathrm{๐ ๐๐}(x,y)=\mathrm{cos}(\stackrel{}{\mathrm{PSP}(x)},\stackrel{}{\mathrm{PSP}(y)}).`$ <sup>8</sup><sup>8</sup>8While admittedly we initially chose this function because it was convenient to work with cosines, post hoc analysis revealed that the corresponding metric space โstretchedโ certain distances in a useful way.
But before proceeding, we note that it is possible that similarity information might yield no extra benefit at all. For instance, we donโt need it if we can reliably identify each class just from some set of distinguishing terms. If we define such terms as frequent ones ($`n20`$) that appear in a single class 50% or more of the time, then we do find many instances; some examples for one author are: โmeaninglessโ, โdisgustingโ (class 0); โpleasantโ, โunevenโ (class 1); and โoscarโ, โgemโ (class 2) for the three-class case, and, in the four-class case, โflatโ, โtediousโ (class 1) versus โstraightforwardโ, โlikeableโ (class 2). Some unexpected distinguishing terms for this author are โlionโ for class 2 (three-class case), and for class 2 in the four-class case, โjenniferโ, for a wide variety of Jennifers.
## 5 Evaluation
This section compares the accuracies of the approaches outlined in Section 3 on the four corpora comprising our scale dataset. (Results using $`L_1`$ error were qualitatively similar.) Throughout, when we refer to something as โsignificantโ, we mean statistically so with respect to the paired $`t`$-test, $`p<.05`$.
The results that follow are based on $`\mathrm{SVM}^{light}`$โs default parameter settings for SVM regression and OVA. Preliminary analysis of the effect of varying the regression parameter $`\epsilon `$ in the four-class case revealed that the default value was often optimal.
The notation โA$`+`$Bโ denotes metric labeling where method A provides the initial label preference function $`\pi `$ and B serves as similarity measure. To train, we first select the meta-parameters $`k`$ and $`\alpha `$ by running 9-fold cross-validation within the training set. Fixing $`k`$ and $`\alpha `$ to those values yielding the best performance, we then re-train A (but with SVM parameters fixed, as described above) on the whole training set. At test time, the nearest neighbors of each item are also taken from the full training set.
### 5.1 Main comparison
Figure 2 summarizes our average 10-fold cross-validation accuracy results. We first observe from the plots that all the algorithms described in Section 3 always definitively outperform the simple baseline of predicting the majority class, although the improvements are smaller in the four-class case. Incidentally, the data was distributed in such a way that the absolute performance of the baseline itself does not change much between the three- and four-class case (which implies that the three-class datasets were relatively more balanced); and Author cโs datasets seem noticeably easier than the others.
We now examine the effect of implicitly using label and item similarity. In the four-class case, regression performed better than OVA (significantly so for two authors, as shown in the righthand table); but for the three-category task, OVA significantly outperforms regression for all four authors. One might initially interprete this โflipโ as showing that in the four-class scenario, item and label similarities provide a richer source of information relative to class-specific characteristics, especially since for the non-majority classes there is less data available; whereas in the three-class setting the categories are better modeled as quite distinct entities.
However, the three-class results for metric labeling on top of OVA and regression (shown in Figure 2 by black versions of the corresponding icons) show that employing explicit similarities always improves results, often to a significant degree, and yields the best overall accuracies. Thus, we can in fact effectively exploit similarities in the three-class case. Additionally, in both the three- and four- class scenarios, metric labeling often brings the performance of the weaker base method up to that of the stronger one (as indicated by the โdisappearanceโ of upward triangles in corresponding table rows), and never hurts performance significantly.
In the four-class case, metric labeling and regression seem roughly equivalent. One possible interpretation is that the relevant structure of the problem is already captured by linear regression (and perhaps a different kernel for regression would have improved its three-class performance). However, according to additional experiments we ran in the four-class situation, the test-set-optimal parameter settings for metric labeling would have produced significant improvements, indicating there may be greater potential for our framework. At any rate, we view the fact that metric labeling performed quite well for both rating scales as a definitely positive result.
### 5.2 Further discussion
Q: Metric labeling looks like itโs just combining SVMs with nearest neighbors, and classifier combination often improves performance. Couldnโt we get the same kind of results by combining SVMs with any other reasonable method?
A: No. For example, if we take the strongest base SVM method for initial label preferences, but replace PSP with the term-overlap-based cosine (TO(cos)), performance often drops significantly. This result, which is in accordance with Section 4โs data, suggests that choosing an item similarity function that correlates well with label similarity is important. (ova$`+`$PSP $``$$``$$``$$``$ ova$`+`$TO(cos) \[3c\]; reg$`+`$PSP $``$ reg$`+`$TO(cos) \[4c\])
Q: Could you explain that notation, please?
A: Triangles point toward the significantly better algorithm for some dataset. For instance, โM $``$$``$$``$ N \[3c\]โ means, โIn the 3-class task, method M is significantly better than N for two author datasets and significantly worse for one dataset (so the algorithms were statistically indistinguishable on the remaining dataset)โ. When the algorithms being compared are statistically indistinguishable on all four datasets (the โno trianglesโ case), we indicate this with an equals sign (โ=โ).
Q: Thanks. Doesnโt Figure 1 show that the positive-sentence percentage would be a good classifier even in isolation, so metric labeling isnโt necessary?
A: No. Predicting class labels directly from the PSP value via trained thresholds isnโt as effective (ova$`+`$PSP $``$$``$$``$$``$ threshold PSP \[3c\]; reg$`+`$PSP $``$$``$ threshold PSP \[4c\]).
Alternatively, we could use only the PSP component of metric labeling by setting the label preference function to the constant function 0, but even with test-set-optimal parameter settings, doing so underperforms the trained metric labeling algorithm with access to an initial SVM classifier (ova$`+`$PSP $``$$``$$``$$``$ 0$`+`$$`\mathrm{PSP}^{}`$ \[3c\]; reg$`+`$PSP $``$$``$ 0$`+`$$`\mathrm{PSP}^{}`$ \[4c\]).
Q: What about using PSP as one of the features for input to a standard classifier?
A: Our focus is on investigating the utility of similarity information. In our particular rating-inference setting, it so happens that the basis for our pairwise similarity measure can be incorporated as an item-specific feature, but we view this as a tangential issue. That being said, preliminary experiments show that metric labeling can be better, barely (for test-set-optimal parameter settings for both algorithms: significantly better results for one author, four-class case; statistically indistinguishable otherwise), although one needs to determine an appropriate weight for the PSP feature to get good performance.
Q: You defined the โmetric transformationโ function $`f`$ as the identity function $`f(d)=d`$, imposing greater loss as the distance between labels assigned to two similar items increases. Can you do just as well if you penalize all non-equal label assignments by the same amount, or does the distance between labels really matter?
A: Youโre asking for a comparison to the Potts model, which sets $`f`$ to the function $`\widehat{f}(d)=1`$ if $`d>0`$, $`0`$ otherwise. In the one setting in which there is a significant difference between the two, the Potts model does worse (ova$`+`$PSP $``$ ova$`\widehat{+}`$PSP \[3c\]). Also, employing the Potts model generally leads to fewer significant improvements over a chosen base method (compare Figure 2โs tables with: reg$`\widehat{+}`$PSP $``$ reg \[3c\]; ova$`\widehat{+}`$PSP $``$$``$ ova \[3c\]; ova$`\widehat{+}`$PSP $`=`$ ova \[4c\]; but note that reg$`\widehat{+}`$PSP $``$ reg \[4c\]). We note that optimizing the Potts model in the multi-label case is NP-hard, whereas the optimal metric labeling with the identity metric-transformation function can be efficiently obtained (see Section 3.3).
Q: Your datasets had many labeled reviews and only one author each. Is your work relevant to settings with many authors but very little data for each?
A: As discussed in Section 2, it can be quite difficult to properly calibrate different authorsโ scales, since the same number of โstarsโ even within what is ostensibly the same rating system can mean different things for different authors. But since you ask: we temporarily turned a blind eye to this serious issue, creating a collection of 5394 reviews by 496 authors with at most 80 reviews per author, where we pretended that our rating conversions mapped correctly into a universal rating scheme. Preliminary results on this dataset were actually comparable to the results reported above, although since we are not confident in the class labels themselves, more work is needed to derive a clear analysis of this setting. (Abusing notation, since weโre already playing fast and loose: \[3c\]: baseline 52.4%, reg 61.4%, reg$`+`$PSP 61.5%, ova (65.4%) $``$ ova$`+`$PSP (66.3%); \[4c\]: baseline 38.8%, reg (51.9%) $``$ reg$`+`$PSP (52.7%), ova (53.8%) $``$ ova$`+`$PSP (54.6%))
In future work, it would be interesting to determine author-independent characteristics that can be used on (or suitably adapted to) data for specific authors.
Q: How about trying โ
A: โYes, there are many alternatives. A few that we tested are described in the Appendix, and we propose some others in the next section. We should mention that we have not yet experimented with all-vs.-all (AVA), another standard binary-to-multi-category classifier conversion method, because we wished to focus on the effect of omitting pairwise information. In independent work on 3-category rating inference for a different corpus, ?) found that regression outperformed AVA, and ?) argue that in principle OVA should do just as well as AVA. But we plan to try it out.
## 6 Related work and future directions
In this paper, we addressed the rating-inference problem, showing the utility of employing label similarity and (appropriate choice of) item similarity โ either implicitly, through regression, or explicitly and often more effectively, through metric labeling.
In the future, we would like to apply our methods to other scale-based classification problems, and explore alternative methods. Clearly, varying the kernel in SVM regression might yield better results. Another choice is ordinal regression , which only considers the ordering on labels, rather than any explicit distances between them; this approach could work well if a good metric on labels is lacking. Also, one could use mixture models (e.g., combine โpositiveโ and โnegativeโ language models) to capture class relationships .
We are also interested in framing multi-class but non-scale-based categorization problems as metric labeling tasks. For example, positive vs. negative vs. neutral sentiment distinctions are sometimes considered in which neutral means either objective or a conflation of objective with a rating of mediocre . (Koppel and Schler in independent work also discuss various types of neutrality.) In either case, we could apply a metric in which positive and negative are closer to objective (or objective+mediocre) than to each other. As another example, hierarchical label relationships can be easily encoded in a label metric.
Finally, as mentioned in Section 3.3, we would like to address the transductive setting, in which one has a small amount of labeled data and uses relationships between unlabeled items, since it is particularly well-suited to the metric-labeling approach and may be quite important in practice.
#### Acknowledgments
We thank Paul Bennett, Dave Blei, Claire Cardie, Shimon Edelman, Thorsten Joachims, Jon Kleinberg, Oren Kurland, John Lafferty, Guy Lebanon, Pradeep Ravikumar, Jerry Zhu, and the anonymous reviewers for many very useful comments and discussion. We learned of Moshe Koppel and Jonathan Schlerโs work while preparing the camera-ready version of this paper; we thank them for so quickly answering our request for a pre-print. Our descriptions of their work are based on that pre-print; we apologize in advance for any inaccuracies in our descriptions that result from changes between their pre-print and their final version. We also thank CMU for its hospitality during the year. This paper is based upon work supported in part by the National Science Foundation (NSF) under grant no. IIS-0329064 and CCR-0122581; SRI International under subcontract no. 03-000211 on their project funded by the Department of the Interiorโs National Business Center; and by an Alfred P. Sloan Research Fellowship. Any opinions, findings, and conclusions or recommendations expressed are those of the authors and do not necessarily reflect the views or official policies, either expressed or implied, of any sponsoring institutions, the U.S. government, or any other entity.
## Appendix A Appendix: other variations attempted
### A.1 Discretizing binary classification
In our setting, we can also incorporate class relations by directly altering the output of a binary classifier, as follows. We first train a standard SVM, treating ratings greater than 0.5 as positive labels and others as negative labels. If we then consider the resulting classifier to output a positivity-preference function $`\pi _+(x)`$, we can then learn a series of thresholds to convert this value into the desired label set, under the assumption that the bigger $`\pi _+(x)`$ is, the more positive the review.<sup>9</sup><sup>9</sup>9 This is not necessarily true: if the classifierโs goal is to optimize binary classification error, its major concern is to increase confidence in the positive/negative distinction, which may not correspond to higher confidence in separating โfive starsโ from โfour starsโ. This algorithm always outperforms the majority-class baseline, but not to the degree that the best of SVM OVA and SVM regression does. ?) independently found in a three-class study that thresholding a positive/negative classifier trained only on clearly positive or clearly negative examples did not yield large improvements.
### A.2 Discretizing regression
In our experiments with SVM regression, we discretized regression output via a set of fixed decision thresholds $`\{0.5,1.5,2.5,\mathrm{}\}`$ to map it into our set of class labels. Alternatively, we can learn the thresholds instead. Neither option clearly outperforms the other in the four-class case. In the three-class setting, the learned version provides noticeably better performance in two of the four datasets. But these results taken together still mean that in many cases, the difference is negligible, and if we had started down this path, we would have needed to consider similar tweaks for one-vs-all SVM as well. We therefore stuck with the simpler version in order to maintain focus on the central issues at hand.
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# Pressure-tuned First-order Phase Transition and Accompanying Resistivity Anomaly in CeZn1-ฮดSb2
## Abstract
The Kondo lattice system CeZn<sub>0.66</sub>Sb<sub>2</sub> is studied by electrical resistivity and ac magnetic susceptibility measurements at several pressures. At $`P=0`$ kbar, ferromagnetic and antiferromagnetic transitions appear at 3.6 and 0.8 K, respectively. The electrical resistivity at $`T_N`$ dramatically changes from the Fisher-Langer type (ferromagnetic like) to the Suezaki-Mori type near $`17`$ kbar, i.e., from a positive divergence to a negative divergence in the temperature derivative of the resistivity. The pressure-induced SM type anomaly, which shows thermal hysteresis, is easily suppressed by small magnetic field (1.9 kOe for 19.8 kbar), indicating a weakly first-order nature of the transition. By subtracting a low-pressure data set, we directly compare the resistivity anomaly with the SM theory without any assumption on backgrounds, where the negative divergence in $`d\rho /dT`$ is ascribed to enhanced critical fluctuations in the presence of superzone gaps.
Rare-earth and actinide compounds with localized $`f`$-electrons at high temperatures show non-magnetic or even a superconducting ground state, which cannot be explained by Hillโs observation that magnetic ordering depends on the $`f`$-atom spacing: magnetism at large $`ff`$ spacing and paramagnetism or superconductivity at small separation.Hill (1970) The importance of $`sf`$ exchange coupling $`J`$ between localized $`f`$-electrons and conduction band electrons for determining the ground state, which was neglected by Hill, is typified in heavy-electron, Kondo-lattice, or mixed-valence systems, where a paramagnetic or a superconducting ground state is observed even though they have large $`ff`$ spacing.Fisk et al. (1988); Thompson (1992) Depending on the hybridization strength $`|JN_F|`$, the ground state is determined through competition between Kondo and Ruderman-Kittel-Kasuya-Yosida (RKKY) interactions, where $`N_F`$ is the conduction band density-of-states at the Fermi energy $`E_F`$. For $`|JN_F|<<1`$, a magnetic state is stabilized because the RKKY interaction that provides coupling between local moments depends geometrically on $`|J|`$ \- $`T_{RKKY}J^2N_F`$, whereas a Kondo singlet is preferred for large $`JN_F`$ because the Kondo interaction that screens local moments depends exponentially on $`|J|`$ \- $`T_K\text{exp}(1/|J|N_F)`$. Doniach (1977)
Another important consequence of the exchange interaction between the localized and conduction electrons is a change in the electronic structure. When an antiferromagnetic structure (AF) with a period incommensurate with the ionic lattice appears, the magnetic superlattice may distort the Fermi surface dramatically,Mackintosh (1962) forming magnetic superzone gaps below $`T_N`$ when the ordering wavevector $`K_A`$ connects portions of the Fermi surface.Salamon (1972) Suezaki and Mori (SM) showed that, when combined with enhanced spin scattering in a $`K=K_A`$ mode in antiferromagnets, this sharp band gap gives rise to a sharp increase in the resistivity or a negative divergence in the temperature derivative near $`T_N`$.Suezaki and Mori (1969) Early measurements on rare-earth metals and order-disorder systems revealed a similar resistivity anomaly, but quantitative analysis has been limited due to the smearing of the gaps by thermal phonons and temperature dependent backgrounds.Craven and Parks (1973); Thomas et al. (1973) Here we report a pressure-tuned first-order phase transition and an accompanying negative divergence in the temperature derivative of the resistivity of CeZn<sub>0.66</sub>Sb<sub>2</sub>. The first-order, SM type transition at $`T_N`$ only appears at intermediate pressures ($`17.3P<25.5`$ kbar), while the transition shows a Fisher-Langer type anomaly, i.e., a positive divergence in $`d\rho /dT`$ in the low pressure limit and a slight slope change at $`P25.5`$ kbar. The first-order anomaly is suppressed by a magnetic field as small as 1.9 kOe at 19.8 kbar, indicating the transition is very weakly first order. By subtracting a low-pressure data set, we directly compare the resistivity anomaly with the SM theory without any assumption on backgrounds, where the negative divergence in $`d\rho /dT`$ is ascribed to enhanced critical fluctuations in the presence of superzone gaps.
CeZn<sub>0.66</sub>Sb<sub>2</sub> was grown with Sb self flux in an evacuated and sealed quartz ampule and crystallizes in the tetragonal ZrCuSi<sub>2</sub> structure with space group P4/nmm. Lee (2005) The electrical resistivity of CeZn<sub>0.66</sub>Sb<sub>2</sub> at 0.3 K is $`4.9\mu \mathrm{\Omega }`$cm and the resistivity ratio is $`\rho (=300\text{K})/\rho (=0.3\text{(}K))14`$. Hydrostatic pressure up to 25 kbar was achieved by using a hybrid Be-Cu/NiCrAl clamp-type pressure cell with silicon fluid as a transmitting medium. At higher pressures, a profiled toroidal anvil clamped device was used with anvils supplied with a boron-epoxy gasket and a teflon capsule filled with glycerol/water mixture with volume ratio 3:2. Superconducting transition temperatures of Tin and Lead were used to determine pressure for the clamp-type and the toroidal anvil cell, respectively. The width of the superconducting transition is independent of pressure and is less than 10 mK up to 55 kbar, indicating that measurements were performed in hydrostatic conditions. Electrical resistivity $`\rho `$ was measured by a standard four-point method with a LR-700 ac resistance bridge (Linear Research) for current flowing perpendicular to the c-axis of CeZn<sub>0.66</sub>Sb<sub>2</sub> ($`I`$ c-axis). AC magnetic susceptibility $`\chi _{\text{ac}}`$ in the plane was measured at $`f=157`$ Hz by a conventional method using primary and secondary pick-up coils mounted inside the pressure cell.
Figure 1 shows the resistivity of CeZn<sub>0.66</sub>Sb<sub>2</sub> as a function of temperature at several pressures. At ambient pressure, a sharp decrease occurs at 3.6 K, corresponding to a ferromagnetic phase transition observed in specific heat measurements. Lee (2005) With increasing pressure, the ferromagnetic transition temperature ($`T_c`$) increases at a rate $`dT_c/dP0.05`$ K/kbar and the resistivity above $`T_c`$ slightly increases (see inset to Fig. 1), which could be explained by the enhanced hybridization between $`f`$ and itinerant electrons with pressure. A slight decrease in $`\rho `$ is also observed at 0.85 K and ambient pressure, which corresponds to an antiferromagnetic transition observed in Ref. . The entropy recovered at $`T_N`$ is about 20 % of $`Rln2`$,Lee (2005) suggesting that the AF state at low temperature is a bulk property, not due to an impurity phase. The antiferromagnetic transition temperature $`T_N`$ slowly increases with pressure, $`dT_N/dP0.02`$ K/kbar. At 17.3 kbar, a peak-like feature with thermal hysteresis appears: the resistivity sharply increases with decreasing temperature, and then decreases. As shown in Fig. 1, with further increasing pressure, the anomaly becomes more pronounced and the transition temperature increases at a much faster rate, $`dT_N/dP0.14`$ K/kbar. Above 25.5 kbar, however, the peak disappears and only a slight slope change occurs at $`T_N`$. A similar pressure-induced resistivity peak was reported in CeRhGe. Asai et al. (2003)
Figure 2 shows representative resistivity data at 19.3 kbar as a function of temperature for several fields. Even though a different piece was used for this measurement, the resistivity anomaly is still reproducible, indicating it is intrinsic. The resistivity peak and thermal hysteresis are steeply suppressed with magnetic field. At as small a magnetic field as 1.9 kOe within the ab-plane, it is totally depressed, suggesting that the AF phase transition is weakly first order. The inset to Fig. 2 gives the temperature derivative of the resistivity. When there is a negative divergence in $`d\rho /dT`$, we assigned the divergent point as the transition temperature $`(T_N)`$, while the maximum point was assigned to $`T_N`$ for the field without negative divergence. $`T_N`$ decreases with magnetic field, $`dT_N/dH0.13\pm 0.04`$ K/kOe.
Figure 3(a) shows the evolution of the temperature derivative of the resistivity for several pressures. At the ferromagnetic transition temperature, $`d\rho /dT`$ diverges for $`TT_c^+`$. We fit the positive divergence to the following form that is commonly used for critical fluctuations for $`T>T_c`$:
$$d\rho /dT=(A/ฯต)(1+|t|^ฯต)+B,$$
(1)
where $`t=(TT_c)/T_c`$. When $`ฯต`$ approaches zero, the above form suggests a logarithmic singularity at $`T_c`$. The inset to Fig. 3(a) magnifies $`d\rho /dT`$ for 19.8 kbar near $`T_c`$ and the solid line is the least-square fit to Eq. (1). The best result was obtained with $`T_c=4.419`$ K and $`ฯต=0.04`$, where the critical exponent is similar to other ferromagnets.Zumsteg and Parks (1970) However, as pointed out by Kadanoff et al.,Kadanoff et al. (1967) the determination of the critical exponent depends on the range of the fit and availability of a number of data points near $`T_c`$. Nevertheless, a similar sharp peak, like that shown in the inset, is observed in the specific heat, Lee (2005) which could be explained by the Fisher-Langer prediction that the magnetic contributions to $`d\rho /dT`$ and the specific heat of a ferromagnet should be proportional because short-range spin-correlations dominate in the temperature dependence of both quantities.Fisher and Langer (1968) Below $`T_c`$, the AF magnetic transition makes it difficult to analyze the critical behavior. Above 25.5 kbar where the resistivity anomaly disappears, the peak in $`d\rho /dT`$ becomes broadened, making a quantitative fit impossible.
The negative divergence in $`d\rho /dT`$ at the N$`\stackrel{ยด}{\text{e}}`$el temperature for intermediate pressures (Fig. 3a) can be understood in terms of combined effects of AF critical fluctuations and superzone gaps below $`T_N`$. In electrical resistivity $`\rho =m/e^2\tau n_{eff}`$, the effective number of carriers $`n_{eff}`$ depends on superzone gaps arising from the additional magnetic lattice periodicity, while the scattering rate $`1/\tau `$ is related to critical scattering of conduction electrons by localized spins. In ferromagnets, spin fluctuations with wave vectors close to $`K_A=0`$ contribute to small angle scattering of conduction electrons, leading to a weak anomaly, i.e., Fisher-Langer type. In antiferromagnets including those with a helical structure, on the other hand, critical fluctuations around the ordering wave vector $`K_A=Q`$ contribute to large angle scattering, leading to a large anomaly. Suezaki and Mori (SM) took into account this critical scattering and predicted the following form:Suezaki and Mori (1969)
$$\begin{array}{cc}d\rho /dT=Bt^{1/3},(T>T_N),\hfill & \\ =B^{^{}}t^{1/3}B_gt^{2/3},(T<T_N).\hfill & \end{array}$$
(2)
For $`T<T_N`$, the first term is due to critical fluctuations and the second term is from long-range order, while only critical fluctuations contribute to the resistivity for $`T>T_N`$. Direct comparison of the critical phenomena between experiments and theory has been limited due to other contributions to the temperature dependent resistivity, such as lattice vibrations. In order to account for other contributions, we subtracted a low-pressure data set (14.2 kbar) because $`T_c`$ is close to that at intermediate pressures even though $`T_N`$ is below 1 K. Figure 3(b) shows representative $`\mathrm{\Delta }|d\rho /dT|`$ for $`T>T_N`$ as a function of $`t(=T/T_N1)`$ for intermediate pressures, where $`\mathrm{\Delta }|d\rho /dT|=|d\rho /dT(P)d\rho /dT(14.2\text{kbar})|`$ and $`T_N`$ is assigned as the negative peak in $`d\rho /dT`$. $`\mathrm{\Delta }|d\rho /dT|`$ for intermediate pressures shows scaling behavior, indicating the validity of the background subtraction. For direct comparison with the SM theory, a power-law form, $`\mathrm{\Delta }d\rho /dT=Bt^\alpha `$, was used for $`T>T_N`$ and the best result was obtained with $`B=0.25`$ and $`\alpha =0.48\pm 0.08`$ by least-squares technique (dashed line in Fig. 3b). The obtained exponent is compatible with the predicted value 1/3 from the SM theory. For $`T<T_N`$, all data sets, similar to those for $`T>T_N`$, collapse on top of each other and Eq. (2) gives a good description of the data with $`B^{^{}}=0.37`$ and $`B_g=0.32`$ (see inset to Fig. 3b). We note that our analysis suffers from the limited temperature range of fitting because the low AF transition temperature makes it difficult to access a reasonable reduced temperature range near $`T_N`$: $`t=0.01`$ for $`\mathrm{\Delta }T=TT_N10`$ mK. As in Cr where a larger exponent was obtained,Akiba and Mitsui (1972) the weakly first order nature of the transition can also complicate the analysis. The relatively good agreement between the experimental data and the SM theory in CeZn<sub>0.66</sub>Sb<sub>2</sub>, however, suggests that the resistivity anomaly mainly comes from critical fluctuations and magnetic gaps , which is consistent with the conclusion from the magnetic field dependence of the anomaly that the resistive transition is weakly first order.
Figure 4(a) shows ac magnetic susceptibility as a function of temperature at several pressures for $`H_{ac}`$c-axis. A sharp peak corresponding to the ferromagnetic transition occurs at 3.5 K for $`P=0`$ kbar and moves toward higher temperature with $`P`$ at the same rate as that of $`T_c`$ determined by the resistivity, as shown in Fig. 4(c). The resistivity anomaly at the N$`\stackrel{ยด}{e}`$el temperature dramatically changes from Fisher-Langer (FL) or ferromagnetic like to Suezaki-Mori (SM) behavior with increasing pressure, i.e., from positive divergence to negative divergence in $`d\rho /dT`$ at $`T_N`$. Salamon claimed that it is prerequisite for magnetic superzone gaps that the ordering wavevector $`K_A`$ should connect portions of the Fermi surface. Salamon (1972) In rare earth metals, spiral spin structures lead to small values of $`K_A`$ which span the Fermi surface, at least in some directions and, therefore, a SM behavior in resistivity. In beta brass where disorder makes it difficult for $`K_A`$ to span the Fermi surface, in contrast, FL behavior was reported.Simons and Salamon (1971) The inset to Fig. 4(a) magnifies the temperature range near $`T_N`$. At low pressures (AF1 in Fig. 4c), there is no clear signature in $`\chi _{ac}`$ corresponding to the AF transition below 1 K (not shown). Above 14.2 kbar (AF2 in Fig. 4c), a peak appears at the temperature corresponding to $`T_N`$ determined from $`\rho `$ and becomes enhanced with pressure (see inset of Fig. 4a and Fig. 4b), implying that the AF structure at intermediate pressures is not simple, but rather has a canted or helical structure. The concurrence of the SM behavior in the resistivity anomaly and the peak feature in $`\chi _{\text{ac}}`$ near 17 kbar is similar to the $`\gamma `$-phase Fe-Mn alloys, Fe<sub>x</sub>Mn<sub>1-x</sub>, where the resistivity anomaly change from FL type to SM type near $`x=0.3`$ was associated with a spin structure change from colinear to non-colinear one Endoh and Ishikawa (1971) and gap formation.Asano and Yamashita (1971) Even though we need to determine exact spin structures from other measurements, such as neutron scattering under pressure, the above analogy suggests that it is only for intermediate pressures where the conditions for magnetic superzone gaps are met in CeZn<sub>0.66</sub>Sb<sub>2</sub>, thus leading to the SM-type resistivity anomaly. Magnetic field dependence of the anomaly is also consistent with the analysis in that the superzone gaps formed at $`H=0`$ Oe are destroyed with magnetic field at intermeidate pressures, thus leading to a change from the SM-type to the FL-type (see Fig. 2).
We have reported a pressure-induced, first-order resistivity anomaly where the resistivity of CeZn<sub>0.66</sub>Sb<sub>2</sub> increases with decreasing temperature and shows thermal hysteresis at the antiferromangetic transition temperature $`T_N`$. By subtracting a low-pressure data set, we directly compared our experiments to Suezaki-Mori theory without any assumption on backgrounds and found reasonably good agreement both below and above $`T_N`$. The dramatic pressure effects on the resistivity anomaly, from the low-pressure Fisher-Langer type to the intermediate-pressure Suezaki-Mori type, are explained in terms of gap formation on the Fermi surface when the AF ordering wavevector $`K_A`$ is tuned to span the Fermi surface by pressure. Magnetic field dependence of the anomaly is also consistent with gap formation at intermediate pressures.
Work at Los Alamos was performed under the auspices of the U.S. Department of Energy/Office of Science. Work at UC-Davis was supported by by NSF Grant No. DMR-0433560. V.A.S. acknowledges the support of Russian Foundation for Basic Research (Grant No. 03-02-17119) and Program โPhysics and Mechanics of Strongly Compressed Matter of Presidium of Russian Academy of Sciencesโ. We benefited from stimulating discussion with M. B. Salamon.
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# IIB Supergravity Revisited
## 1 Introduction
IIB supergravity is the low energy effective action of type-IIB superstring theory. Its scalar sector describes the coset manifold $`SL(2,)/SO(2)SU(1,1)/U(1)`$, whose isometry $`SL(2,)`$ is a symmetry of the low energy theory. Since the isometry acts non-trivially on the dilaton, the full perturbative string theory does not preserve the symmetry, but the conjecture is that non-perturbatively an $`SL(2,)`$ subgroup of the full symmetry group of the low energy action survives .
The particular feature of type-IIB string theory with respect to the other theories of closed oriented strings is that it is symmetric under the orientation reversal of the fundamental string. Ten-dimensional type-I string theory is obtained from type-IIB through an orientifold projection that gauges this symmetry, and tadpole cancellation requires the introduction of an open sector, corresponding to D9-branes. The standard supersymmetric projection gives rise to the type-I superstring, with gauge group $`SO(32)`$ , while a non-supersymmetric, anomaly-free projection gives rise to a model with gauge group $`USp(32)`$ , in which supersymmetry is realized on the bulk and spontaneously broken on the branes .
In the low-energy effective action, the closed sector of type-I strings is obtained by performing a consistent $`_2`$ truncation of the IIB supergravity, while the open sector corresponds to the first order in the low-energy expansion of the D9-brane action in a type-I background. In it was shown that the $`_2`$ symmetry responsible for this truncation can be performed in two ways, and in a flat background, with all bulk fields put to zero, the D9-brane action reduces in one case to the Volkov-Akulov action , and in the other case to a constant. In these results were extended to a generic background, showing that also in the curved case there are two possibilities of performing the truncation. In one case one gets a dilaton tadpole and a RR tadpole plus goldstino couplings, which is basically the one-brane equivalent of the Sugimoto model, while in the other case the goldstino couplings vanish and one is left with a dilaton and a RR tadpole, which is the one-brane equivalent of the supersymmetric model. In order to truncate the theory in the brane sector, the โdemocratic formulationโ of IIB supergravity was derived . This amounts to an extension of the supersymmetry algebra, so that both the RR fields and their magnetic duals appear on the same footing. The closure of the algebra then requires the field strengths of these fields to be related by duality conditions. The result is that, together with the RR forms $`C^{(2n)}`$, $`n=0,\mathrm{},4`$ associated with D-branes of non-vanishing codimension, the algebra naturally includes a RR ten-form $`C^{(10)}`$, with respect to which the spacetime-filling D9-branes are electrically charged. This field does not have any field strength, and correspondingly an object charged with respect to it can be consistently included in the theory only when one performs a type-I truncation, so that the resulting overall RR charge vanishes. The analysis of also showed that an additional ten-form $`B^{(10)}`$ can be introduced in the algebra, and this form survives a different $`_2`$ truncation, projecting out all the RR-fields. In the string frame, the tension of a spacetime-filling brane electrically charged with respect to $`B^{(10)}`$ would scale like $`g_S^2`$, instead of $`g_S^4`$, thus implying that the brane action for this object can not be obtained performing an $`S`$-duality transformation on the D9-brane effective action . We are therefore facing a problem, since two ten-forms are known in IIB supergravity, but they do not form a doublet with respect to $`SL(2,)`$.
In this paper we will clarify this issue. We want to obtain all the possible independent ten-forms that can be added to 10-dimensional IIB supergravity, with their assignment to representations of $`SL(2,)`$. In order to perform this analysis, we express the theory in a โ$`SU(1,1)`$-democratic formulationโ, in which all the forms, not only the RR ones, and their magnetic duals are described in a $`SU(1,1)`$-covariant way. We use the notation of , so that the scalars parametrize the coset $`SU(1,1)/U(1)`$, while the two two-forms, as well as their duals, form a doublet of $`SU(1,1)`$. The eight-forms, dual to the scalars, transform as a triplet of $`SU(1,1)`$, with the field strengths satisfying an $`SU(1,1)`$ invariant constraint . Eventually, we find that the algebra includes a doublet and a quadruplet of ten-forms <sup>1</sup><sup>1</sup>1Gauge fields of maximal rank have been explored in the literature ., and the dilaton dependence of the supersymmetry transformation of these objects shows that the RR ten-form belongs to the quadruplet. We claim that no other independent ten-forms can be added to the algebra. In summary, we find the following bosonic field content:
$$e_\mu ^a,V_+^\alpha ,V_{}^\alpha ,A_{(2)}^\alpha ,A_{(4)},A_{(6)}^\alpha ,A_{(8)}^{(\alpha \beta )},A_{(10)}^\alpha ,A_{(10)}^{(\alpha \beta \gamma )},$$
(1)
where $`e_\mu ^a`$ is the zehnbein, $`(V_+^\alpha ,V_{}^\alpha )`$ parametrizes the $`SU(1,1)/U(1)`$ coset, $`\alpha =1,2`$ is an $`SU(1,1)`$ index and the subindex $`(n)`$ indicates the rank of the potential.
This paper will be devoted to the construction and the properties of the extended IIB supergravity theory (1). Clearly the properties of the dual forms and ten-forms have implications for the structure of the brane spectrum, dualities, etc. These aspects of this work will be addressed in a forthcoming paper .
The structure of the paper is as follows. The main result, the supersymmetry transformation rules and algebra of the extended IIB-supergravity theory in the $`SU(1,1)/U(1)`$ formulation, are given in section 5. In section 6 these results are rewritten in a $`U(1)`$ gauge in the Einstein frame and in the string frame. In this section we also recover the Ramond-Ramond โharmonicaโ of and then extend it to the Neveu-Schwarz forms. We also list the action of $`S`$-duality on all form fields. The preceding sections lead up to these results and sketch the derivation. In section 2 we review the $`SU(1,1)`$-covariant notation of . In section 3 we introduce in the algebra the six- and the eight-forms dual to the two-forms and the scalars respectively. Section 4 contains the analysis of the ten-forms. We finally conclude with a summary of our results and a discussion. Some basic formulas and truncations to $`N=1`$ supergravity can be found in the Appendices.
## 2 The $`SU(1,1)`$-covariant formulation
In this section we review the notation and the results of .
The theory contains the graviton, two scalars, two two-forms and a self-dual four-form in the bosonic sector, together with a complex left-handed gravitino and a complex right-handed spinor in the fermionic sector. We will use the mostly-minus spacetime signature convention throughout the paper. The two scalars parametrize the coset $`SU(1,1)/U(1)`$, that can be described in terms of the $`SU(1,1)`$ matrix ($`\alpha ,\beta =1,2`$)
$$U=(V_{}^\alpha V_+^\alpha ),$$
(2)
satisfying the constraint
$$V_{}^\alpha V_+^\beta V_+^\alpha V_{}^\beta =ฯต^{\alpha \beta },$$
(3)
with $`(V_{}^1)^{}=V_+^2`$, where $`\alpha =1,2`$ is an $`SU(1,1)`$ index and $`+`$ and $``$ denote the $`U(1)`$ charge, and $`ฯต^{12}=ฯต_{12}=1`$. From the left-invariant 1-form
$$U^1_\mu U=\left(\begin{array}{cc}iQ_\mu & P_\mu \\ P_\mu ^{}& iQ_\mu \end{array}\right)$$
(4)
one reads off the $`U(1)`$-covariant quantity
$$P_\mu =ฯต_{\alpha \beta }V_+^\alpha _\mu V_+^\beta ,$$
(5)
that has charge 2, and the $`U(1)`$ connection
$$Q_\mu =iฯต_{\alpha \beta }V_{}^\alpha _\mu V_+^\beta .$$
(6)
Note that
$`P_\mu V_{}^\alpha `$ $`=`$ $`D_\mu V_+^\alpha ,`$ (7)
$`P_\mu ^{}V_+^\alpha `$ $`=`$ $`D_\mu V_{}^\alpha ,`$ (8)
where the derivative $`D`$ is covariant with respect to $`U(1)`$. The two-forms are collected in an $`SU(1,1)`$ doublet $`A_{\mu \nu }^\alpha `$ satisfying the constraint
$$(A_{\mu \nu }^1)^{}=A_{\mu \nu }^2.$$
(9)
The corresponding field strengths
$$F_{\mu \nu \rho }^\alpha =3_{[\mu }A_{\nu \rho ]}^\alpha $$
(10)
are invariant with respect to the gauge transformations
$$\delta A_{\mu \nu }^\alpha =2_{[\mu }\mathrm{\Lambda }_{\nu ]}^\alpha .$$
(11)
The four-form is invariant under $`SU(1,1)`$, and varies as
$$\delta A_{\mu \nu \rho \sigma }=4_{[\mu }\mathrm{\Lambda }_{\nu \rho \sigma ]}\frac{i}{4}ฯต_{\alpha \beta }\mathrm{\Lambda }_{[\mu }^\alpha F_{\nu \rho \sigma ]}^\beta $$
(12)
under four-form and two-form gauge transformations, so that the gauge-invariant five-form field-strength is
$$F_{\mu \nu \rho \sigma \tau }=5_{[\mu }A_{\nu \rho \sigma \tau ]}+\frac{5i}{8}ฯต_{\alpha \beta }A_{[\mu \nu }^\alpha F_{\rho \sigma \tau ]}^\beta .$$
(13)
This five-form satisfies the self-duality condition
$$F^{\mu _1\mathrm{}\mu _5}=\frac{1}{5!}ฯต^{\mu _1\mathrm{}\mu _5\nu _1\mathrm{}\nu _5}F_{\nu _1\mathrm{}\nu _5}.$$
(14)
It is convenient to define the complex three-form
$$G_{\mu \nu \rho }=ฯต_{\alpha \beta }V_+^\alpha F_{\mu \nu \rho }^\beta ,$$
(15)
that is an $`SU(1,1)`$ singlet with $`U(1)`$ charge 1. Finally the gravitino $`\psi _\mu `$ is complex left-handed with $`U(1)`$ charge $`1/2`$, while the spinor $`\lambda `$ is complex right-handed with $`U(1)`$ charge $`3/2`$.
In the field equations for this model were derived by requiring the closure of the supersymmetry algebra. All these equations can be derived from a lagrangian, imposing eq. (14) only after varying <sup>2</sup><sup>2</sup>2A lagrangian formulation for self dual forms has been developed in , and then applied in to the ten-dimensional IIB supergravity. It corresponds to the introduction of an additional scalar auxiliary field, and the self-duality condition results from the gauge fixing (that can not be imposed directly on the action) of additional local symmetries.. It is interesting to study in detail the kinetic term for the scalar fields,
$$_{scalar}=\frac{e}{2}P_\mu ^{}P^\mu .$$
(16)
The complex variable
$$z=\frac{V_{}^2}{V_{}^1}$$
(17)
is invariant under local $`U(1)`$ transformations, and so it is a good coordinate for the scalar manifold. Under the $`SU(1,1)`$ transformation
$$\left(\begin{array}{c}V_{}^1\\ V_{}^2\end{array}\right)\left(\begin{array}{cc}\alpha & \beta \\ \overline{\beta }& \overline{\alpha }\end{array}\right)\left(\begin{array}{c}V_{}^1\\ V_{}^2\end{array}\right),$$
(18)
that is an isometry of the scalar manifold, $`z`$ transforms as
$$z\frac{\overline{\alpha }z+\overline{\beta }}{\beta z+\alpha }.$$
(19)
The variable $`z`$ parametrizes the unit disc, $`|z|<1`$, and the kinetic term assumes the form
$$_{scalar}=\frac{e}{2}\frac{_\mu z^\mu \overline{z}}{(1z\overline{z})^2}.$$
(20)
The further change of variables
$$z=\frac{1+i\tau }{1i\tau }$$
(21)
maps the disc in the complex upper-half plane, $`\mathrm{Im}\tau >0`$, and in terms of $`\tau `$ the transformations (18) become
$$\tau \frac{a\tau +b}{c\tau +d},$$
(22)
where
$$\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)SL(2,),$$
(23)
while the scalar lagrangian takes the form
$$_{scalar}=\frac{e}{8}\frac{_\mu \tau ^\mu \overline{\tau }}{(\mathrm{Im}\tau )^2}.$$
(24)
Expressing $`\tau `$ in terms of the RR scalar and the dilaton,
$$\tau =\mathrm{}+ie^\varphi $$
(25)
and performing the Weyl rescaling $`g_{(E)\mu \nu }e^{\varphi /2}g_{(S)\mu \nu }`$ one ends up with the standard form of the kinetic term of the scalars in IIB supergravity in the string frame.
The supersymmetry transformations that leave the field equations of invariant are
$`\delta e_\mu {}_{}{}^{a}=i\overline{ฯต}\gamma ^a\psi _\mu +i\overline{ฯต}_C\gamma ^a\psi _{\mu C},`$
$`\delta \psi _\mu =D_\mu ฯต+\frac{i}{480}F_{\mu \nu _1\mathrm{}\nu _4}\gamma ^{\nu _1\mathrm{}\nu _4}ฯต+\frac{1}{96}G^{\nu \rho \sigma }\gamma _{\mu \nu \rho \sigma }ฯต_C\frac{3}{32}G_{\mu \nu \rho }\gamma ^{\nu \rho }ฯต_C,`$
$`\delta A_{\mu \nu }^\alpha =V_{}^\alpha \overline{ฯต}\gamma _{\mu \nu }\lambda +V_+^\alpha \overline{ฯต}_C\gamma _{\mu \nu }\lambda _C+4iV_{}^\alpha \overline{ฯต}_C\gamma _{[\mu }\psi _{\nu ]}+4iV_+^\alpha \overline{ฯต}\gamma _{[\mu }\psi _{\nu ]C},`$
$`\delta A_{\mu \nu \rho \sigma }=\overline{ฯต}\gamma _{[\mu \nu \rho }\psi _{\sigma ]}\overline{ฯต}_C\gamma _{[\mu \nu \rho }\psi _{\sigma ]C}\frac{3i}{8}ฯต_{\alpha \beta }A_{[\mu \nu }^\alpha \delta A_{\rho \sigma ]}^\beta ,`$
$`\delta \lambda =iP_\mu \gamma ^\mu ฯต_C\frac{i}{24}G_{\mu \nu \rho }\gamma ^{\mu \nu \rho }ฯต,`$
$`\delta V_+^\alpha =V_{}^\alpha \overline{ฯต}_C\lambda ,`$
$`\delta V_{}^\alpha =V_+^\alpha \overline{ฯต}\lambda _C.`$ (26)
where we denote with $`\mathrm{\Psi }_C`$ the complex (Majorana) conjugate of $`\mathrm{\Psi }`$. The commutator $`[\delta _1,\delta _2]`$ of two supersymmetry transformations of (26) closes on all the local symmetries of the theory, provided one uses the fermionic field equations and the self-duality condition of eq. (14). To lowest order in the fermions, the parameters of the resulting general coordinate transformation, four-form gauge transformation and two-form gauge transformation are<sup>3</sup><sup>3</sup>3We only present the parameters of translations and the two- and four-form gaugetransformations. The parameters of other local symmetries, namely supersymmetry, local Lorentz and local $`U(1)`$ are not used in the analysis of the next sections, and are given in .
$`\xi ^\mu =i\overline{ฯต}_2\gamma ^\mu ฯต_1+i\overline{ฯต}_{2C}\gamma ^\mu ฯต_{1C},`$
$`\mathrm{\Lambda }_\mu ^\alpha =A_{\mu \nu }^\alpha \xi ^\nu 2i[V_+^\alpha \overline{ฯต}_2\gamma _\mu ฯต_{1C}+V_{}^\alpha \overline{ฯต}_{2C}\gamma _\mu ฯต_1],`$
$`\mathrm{\Lambda }_{\mu \nu \rho }=A_{\mu \nu \rho \sigma }\xi ^\sigma {\displaystyle \frac{1}{4}}[\overline{ฯต}_2\gamma _{\mu \nu \rho }ฯต_1\overline{ฯต}_{2C}\gamma _{\mu \nu \rho }ฯต_{1C}]`$
$`\frac{3}{8}ฯต_{\alpha \beta }A_{[\mu \nu }^\alpha \left(V_+^\beta \overline{ฯต}_2\gamma _{\rho ]}ฯต_{1C}+V_{}^\beta \overline{ฯต}_{2C}\gamma _{\rho ]}ฯต_1\right).`$ (27)
In the next section we will extend the algebra in order to include the magnetic duals of the scalars and of the two-form, in such a way that the supersymmetry algebra still closes, once the proper duality relations are used. Once we obtain the supersymmetry transformation of the six- and the eight-forms that are compatible with the algebra obtained from eq. (26), we will include in Section 4 all the possible independent ten-forms that this algebra allows.
## 3 Six-forms and eight-forms
In this section we show how the algebra of eq. (26) is extended introducing the forms magnetically dual to the scalars and the two-forms. As anticipated, closure of the supersymmetry algebra requires the field strengths of these forms to be related to $`P_\mu `$ and the field strengths of the two-forms by suitable duality relations. Generalizing what happens for the four-form (see eqs. (12) and (13)), we will see that the gauge transformations of these fields involve the gauge parameters of all the lower rank forms, and the gauge invariant field strengths will therefore contain lower rank forms as well. After introducing our Ansatz for these field strengths and gauge transformations, the supersymmetry transformations of these fields will then be determined requiring the closure of the supersymmetry algebra. As in the previous section, we will not consider terms higher than quadratic in the fermi fields.
### 3.1 Six-forms
We want to obtain the gauge and supersymmetry transformations for the doublet of six-forms $`A_{\mu _1\mathrm{}\mu _6}^\alpha `$, which are the magnetic duals of the two-forms and thus satisfy the reality condition
$$(A^1)_{\mu _1\mathrm{}\mu _6}^{}=A_{\mu _1\mathrm{}\mu _6}^2.$$
(28)
Generalizing what one obtains for the four-form, we expect the supersymmetry transformation of the six-forms to contain terms involving only spinors and terms containing forms of lower rank. The condition of eq. (28), as well as the requirement that all the terms must have vanishing local $`U(1)`$ charge, fixes the most general transformation of the doublet to be
$`\delta A_{\mu _1\mathrm{}\mu _6}^\alpha `$ $`=`$ $`aV_{}^\alpha \overline{ฯต}\gamma _{\mu _1\mathrm{}\mu _6}\lambda +a^{}V_+^\alpha \overline{ฯต}_C\gamma _{\mu _1\mathrm{}\mu _6}\lambda _C`$ (29)
$`+`$ $`bV_{}^\alpha \overline{ฯต}_C\gamma _{[\mu _1\mathrm{}\mu _5}\psi _{\mu _6]}b^{}V_+^\alpha \overline{ฯต}\gamma _{[\mu _1\mathrm{}\mu _5}\psi _{\mu _6]C}`$
$`+`$ $`cA_{[\mu _1\mathrm{}\mu _4}\delta A_{\mu _5\mu _6]}^\alpha `$
$`+`$ $`d\delta A_{[\mu _1\mathrm{}\mu _4}A_{\mu _5\mu _6]}^\alpha `$
$`+`$ $`ieฯต_{\beta \gamma }\delta A_{[\mu _1\mu _2}^\beta A_{\mu _3\mu _4}^\gamma A_{\mu _5\mu _6]}^\alpha .`$
We want to consider the commutator $`[\delta _1,\delta _2]`$ of two such transformations, to lowest order in the fermi fields.
We first take into account the terms involving the spinors, i.e., the first two lines in eq. (29). Those terms produce the gauge transformation for the six-forms
$`\delta A_{\mu _1\mathrm{}\mu _6}^\alpha `$ $`=`$ $`6_{[\mu _1}\mathrm{\Lambda }_{\mu _2\mathrm{}\mu _6]}^\alpha `$ (30)
$`=`$ $`12i_{[\mu _1}(aV_+^\alpha \overline{ฯต}_2\gamma _{\mu _2\mathrm{}\mu _6]}ฯต_{1C}+a^{}V_{}^\alpha \overline{ฯต}_{2C}\gamma _{\mu _2\mathrm{}\mu _6]}ฯต_1)`$
if the constraint
$$12ia^{}=b$$
(31)
is imposed, while the other terms that are produced are
$`20iaF_{[\mu _1\mu _2\mu _3}^\alpha (\overline{ฯต}_{2C}\gamma _{\mu _4\mu _5\mu _6]}ฯต_{1C}\overline{ฯต}_2\gamma _{\mu _4\mu _5\mu _6]}ฯต_1)`$
$`\frac{1}{6}aฯต_{\mu _1\mathrm{}\mu _6\sigma \mu \nu \rho }S^{\alpha \beta }ฯต_{\beta \gamma }F^{\gamma ;\mu \nu \rho }\xi ^\sigma ,`$ (32)
where we have defined
$$S^{\alpha \beta }=V_{}^\alpha V_+^\beta +V_+^\alpha V_{}^\beta $$
(33)
and we have assumed that $`a`$ is imaginary. Note that $`S^{\alpha \beta }`$ satisfies
$$S^{\alpha \beta }ฯต_{\beta \gamma }S^{\gamma \delta }ฯต_{\delta ฯต}=\delta _ฯต^\alpha .$$
(34)
Observe that there are no terms involving the five-form field strength. Without loss of generality, we fix
$$a=i$$
(35)
from now on. In order for the last term in (32) to produce a general coordinate transformation with the right coefficient as dictated by eq. (27), we impose the duality relation <sup>4</sup><sup>4</sup>4Note that this duality relation induces field equations for the potentials.
$$F_{\mu _1\mathrm{}\mu _7}^\alpha =\frac{i}{3!}ฯต_{\mu _1\mathrm{}\mu _7\mu \nu \rho }S^{\alpha \beta }ฯต_{\beta \gamma }F^{\gamma ;\mu \nu \rho },$$
(36)
where $`F_{\mu _1\mathrm{}\mu _7}^\alpha =7_{[\mu _1}A_{\mu _2\mathrm{}\mu _7]}^\alpha +\mathrm{}`$ are the field strengths of the six-forms, and the dots stand for terms involving lower rank forms that we will determine in the following. Note that the second term of eq. (32) contains, together with a general coordinate transformation, a gauge transformation of parameter
$$\mathrm{\Lambda }^{}{}_{\mu _1\mathrm{}\mu _5}{}^{\alpha }=A_{\mu _1\mathrm{}\mu _5\sigma }^\alpha \xi ^\sigma .$$
(37)
The $`SU(1,1)`$-invariant quantities
$`G_{\mu _1\mathrm{}\mu _7}=ฯต_{\alpha \beta }V_+^\alpha F_{\mu _1\mathrm{}\mu _7}^\beta ,G_{\mu _1\mathrm{}\mu _7}^{}=ฯต_{\alpha \beta }V_{}^\alpha F_{\mu _1\mathrm{}\mu _7}^\beta ,`$ (38)
which have $`U(1)`$ charge $`+1`$ and $`1`$ respectively, satisfy
$`G_{\mu _1\mathrm{}\mu _7}^{(7)}=\frac{i}{3!}ฯต_{\mu _1\mathrm{}\mu _7\mu \nu \rho }G^{\mu \nu \rho },G_{\mu _1\mathrm{}\mu _7}^{}=\frac{i}{3!}ฯต_{\mu _1\mathrm{}\mu _7\mu \nu \rho }G^{\mu \nu \rho }.`$ (39)
In order to proceed further, in analogy with eq. (13) we make the following Ansatz for the seven-form field strengths:
$$F_{\mu _1\mathrm{}\mu _7}^\alpha =7_{[\mu _1}A_{\mu _2\mathrm{}\mu _7]}^\alpha +\alpha A_{[\mu _1\mu _2}^\alpha F_{\mu _3\mathrm{}\mu _7]}+\beta F_{[\mu _1\mathrm{}\mu _3}^\alpha A_{\mu _4\mathrm{}\mu _7]}.$$
(40)
For these forms to be gauge invariant, the must transform non-trivially with respect to the two-form and four-form gauge transformations. The result is
$$\delta A_{\mu _1\mathrm{}\mu _6}^\alpha =\frac{2}{7}\alpha \mathrm{\Lambda }_{[\mu _1}^\alpha F_{\mu _2\mathrm{}\mu _6]}+\frac{4}{7}\beta F_{[\mu _1\mathrm{}\mu _3}^\alpha \mathrm{\Lambda }_{\mu _4\mathrm{}\mu _6]},$$
(41)
and gauge invariance requires
$$\beta =\frac{10}{3}\alpha .$$
(42)
Now we come back to the commutator. The terms that are left are the ones coming from the last three lines in eq. (29), together with the first line in eq. (32) and the terms coming from (40) in the second line of eq. (32). All these terms have to produce gauge transformations according to (41), with parameters given from eqs. (27), possibly together with additional gauge transformations. The end result is that one produces the additional gauge transformations
$`\mathrm{\Lambda }^{}{}_{}{}^{}_{\mu _1\mathrm{}\mu _5}^{\alpha }`$ $`=`$ $`\frac{2i}{3}cA_{[\mu _1\mathrm{}\mu _4}(V_+^\alpha \overline{ฯต}_2\gamma _{\mu _5]}ฯต_{1C}+V_{}^\alpha \overline{ฯต}_{2C}\gamma _{\mu _5]}ฯต_1)`$ (43)
$``$ $`\frac{1}{6}dA_{[\mu _1\mu _2}^\alpha (\overline{ฯต}_2\gamma _{\mu _3\mathrm{}\mu _5]}ฯต_1\overline{ฯต}_{2C}\gamma _{\mu _3\mathrm{}\mu _5]}ฯต_{1C}),`$
while all the coefficients are uniquely determined to be
$$c=40,d=20,e=\frac{15}{2},\alpha =28.$$
(44)
Summarizing, we get that the supersymmetry transformations of the six-forms are
$`\delta A_{\mu _1\mathrm{}\mu _6}^\alpha `$ $`=`$ $`iV_{}^\alpha \overline{ฯต}\gamma _{\mu _1\mathrm{}\mu _6}\lambda iV_+^\alpha \overline{ฯต}_C\gamma _{\mu _1\mathrm{}\mu _6}\lambda _C`$ (45)
$`+`$ $`12V_{}^\alpha \overline{ฯต}_C\gamma _{[\mu _1\mathrm{}\mu _5}\psi _{\mu _6]}12V_+^\alpha \overline{ฯต}\gamma _{[\mu _1\mathrm{}\mu _5}\psi _{\mu _6]C}`$
$`+`$ $`40A_{[\mu _1\mathrm{}\mu _4}\delta A_{\mu _5\mu _6]}^\alpha `$
$``$ $`20\delta A_{[\mu _1\mathrm{}\mu _4}A_{\mu _5\mu _6]}^\alpha `$
$`+`$ $`\frac{15i}{2}ฯต_{\beta \gamma }\delta A_{[\mu _1\mu _2}^\beta A_{\mu _3\mu _4}^\gamma A_{\mu _5\mu _6]}^\alpha .`$
The doublet of seven-form field strengths is
$$F_{\mu _1\mathrm{}\mu _7}^\alpha =7_{[\mu _1}A_{\mu _2\mathrm{}\mu _7]}^\alpha +28A_{[\mu _1\mu _2}^\alpha F_{\mu _3\mathrm{}\mu _7]}\frac{280}{3}F_{[\mu _1\mathrm{}\mu _3}^\alpha A_{\mu _4\mathrm{}\mu _7]}.$$
(46)
This is gauge invariant with respect to the transformations of the two-forms, the four-form and the six-forms, where
$$\delta A_{\mu _1\mathrm{}\mu _6}^\alpha =6_{[\mu _1}\mathrm{\Lambda }_{\mu _2\mathrm{}\mu _6]}^\alpha 8\mathrm{\Lambda }_{[\mu _1}^\alpha F_{\mu _2\mathrm{}\mu _6]}\frac{160}{3}F_{[\mu _1\mathrm{}\mu _3}^\alpha \mathrm{\Lambda }_{\mu _4\mathrm{}\mu _6]}.$$
(47)
Moreover, the six-form gauge transformation parameter resulting from the commutator of two supersymmetry transformations is
$`\mathrm{\Lambda }_{\mu _1\mathrm{}\mu _5}^\alpha `$ $`=`$ $`A_{\mu _1\mathrm{}\mu _5\sigma }^\alpha \xi ^\sigma +2(V_+^\alpha \overline{ฯต}_2\gamma _{\mu _1\mathrm{}\mu _5}ฯต_1V_{}^\alpha \overline{ฯต}_{2C}\gamma _{\mu _1\mathrm{}\mu _5}ฯต_{1C})`$ (48)
$``$ $`\frac{80i}{3}A_{[\mu _1\mathrm{}\mu _4}(V_+^\alpha \overline{ฯต}_2\gamma _{\mu _5]}ฯต_{1C}+V_{}^\alpha \overline{ฯต}_{2C}\gamma _{\mu _5]}ฯต_1)`$
$`+`$ $`\frac{10}{3}A_{[\mu _1\mu _2}^\alpha (\overline{ฯต}_2\gamma _{\mu _3\mathrm{}\mu _5]}ฯต_1\overline{ฯต}_{2C}\gamma _{\mu _3\mathrm{}\mu _5]}ฯต_{1C}),`$
as results from eqs. (30), (37) and (43). Finally, a comment is in order. At first sight, the Ansatz we made for the field strengths in eq. (40) does not seem to be the most general one, since one could in principle include a term of the form $`iฯต_{\beta \gamma }A_{[\mu _1\mu _2}^aA_{\mu _3\mu _4}^\beta F_{\mu _5\mathrm{}\mu _7]}^\gamma `$. The reason why we did not include it is that one can always reabsorb such a term by performing a redefinition of the six-forms of the type $`A_{\mu _1\mathrm{}\mu _6}^\alpha A_{\mu _1\mathrm{}\mu _6}^\alpha +\gamma A_{[\mu _1\mu _2}^\alpha A_{\mu _3\mathrm{}\mu _6]}`$, and choose $`\gamma `$ so that this term vanishes. This freedom will be used to constrain the form of the field strengths of the eight-forms as well, as we will see in the next subsection.
### 3.2 Eight-forms
The eight-forms are the magnetic duals of the scalars. As we reviewed in Section 2, the scalars are described in terms of the left-invariant 1-form of eq. (4), transforming in the adjoint of $`SU(1,1)`$, and propagating two real degrees of freedom because of local $`U(1)`$ invariance. One therefore expects a triplet of eight-forms (as observed in ) <sup>5</sup><sup>5</sup>5A similar observation was made for the curvatures in ., that we denote by $`A_{\mu _1\mathrm{}\mu _8}^{\alpha \beta }`$, symmetric under $`\alpha \beta `$, and satisfying the reality condition
$$(A^{11})_{\mu _1\mathrm{}\mu _8}^{}=A_{\mu _1\mathrm{}\mu _8}^{22},(A^{12})_{\mu _1\mathrm{}\mu _8}^{}=A_{\mu _1\mathrm{}\mu _8}^{12}.$$
(49)
The fact that only two scalars propagate will result in a constraint for the field strengths of these eight-forms . This is exactly what we are going to show in this subsection. Following the same arguments as in the previous subsection, we write the most general supersymmetry transformations for the eight-forms, compatible with the reality condition and with $`U(1)`$ invariance, consisting of terms that only involve the spinors and terms containing the lower rank forms and their supersymmetry transformations. The result is
$`\delta A_{\mu _1\mathrm{}\mu _8}^{\alpha \beta }`$ $`=`$ $`aV_+^\alpha V_+^\beta \overline{ฯต}\gamma _{\mu _1\mathrm{}\mu _8}\lambda _C+a^{}V_{}^\alpha V_{}^\beta \overline{ฯต}_C\gamma _{\mu _1\mathrm{}\mu _8}\lambda `$ (50)
$`+`$ $`bV_+^{(\alpha }V_{}^{\beta )}\overline{ฯต}\gamma _{[\mu _1\mathrm{}\mu _7}\psi _{\mu _8]}b^{}V_+^{(\alpha }V_{}^{\beta )}\overline{ฯต}_C\gamma _{[\mu _1\mathrm{}\mu _7}\psi _{\mu _8]C}`$
$`+`$ $`cA_{[\mu _1\mathrm{}\mu _6}^{(\alpha }\delta A_{\mu _7\mu _8]}^{\beta )}+dA_{[\mu _1\mu _2}^{(\alpha }\delta A_{\mu _3\mathrm{}\mu _8]}^{\beta )}+ieA_{[\mu _1\mu _2}^{(\alpha }A_{\mu _3\mu _4}^{\beta )}ฯต_{\gamma \delta }A_{\mu _5\mu _6}^\gamma \delta A_{\mu _7\mu _8]}^\delta `$
$`+`$ $`fA_{[\mu _1\mu _2}^{(\alpha }A_{\mu _3\mu _4}^{\beta )}\delta A_{\mu _5\mathrm{}\mu _8]}+gA_{[\mu _1\mathrm{}\mu _4}A_{\mu _5\mu _6}^{(\alpha }\delta A_{\mu _7\mu _8]}^{\beta )}.`$
We first consider the contributions coming from the first two lines of eq. (50), in order to get a relation between $`a`$ and $`b`$. We obtain the gauge transformation
$`\delta A_{\mu _1\mathrm{}\mu _8}^{\alpha \beta }`$ $`=`$ $`8_{[\mu _1}\mathrm{\Lambda }_{\mu _2\mathrm{}\mu _8]}^{\alpha \beta }`$ (51)
$`=`$ $`4ia_{[\mu _1}\left[S^{\alpha \beta }(\overline{ฯต}_2\gamma _{\mu _2\mathrm{}\mu _8]}ฯต_1\overline{ฯต}_{2C}\gamma _{\mu _2\mathrm{}\mu _8]}ฯต_{1C})\right]`$
together with the terms
$`28ia(V_+^{(\alpha }\overline{ฯต}_2\gamma _{[\mu _1\mathrm{}\mu _5}ฯต_1V_{}^{(\alpha }\overline{ฯต}_{2C}\gamma _{[\mu _1\mathrm{}\mu _5}ฯต_{1C})F_{\mu _6\mathrm{}\mu _8]}^{\beta )}`$
$`4a(V_+^{(\alpha }\overline{ฯต}_2\gamma _{[\mu _1}ฯต_1+V_{}^{(\alpha }\overline{ฯต}_{2C}\gamma _{[\mu _1}ฯต_{1C})F_{\mu _2\mathrm{}\mu _8]}^{\beta )}`$
$`aฯต_{\mu _1\mathrm{}\mu _8\sigma \tau }\xi ^\sigma (V_+^\alpha V_+^\beta P^\tau V_{}^\alpha V_{}^\beta P^\tau ),`$ (52)
provided that
$$8ia=b$$
(53)
and $`a`$ is chosen to be imaginary. Fixing, without loss of generality,
$$a=i,$$
(54)
one finds that the last term in eq. (52) contains the correct general coordinate transformation, plus an gauge transformation of parameter
$$\mathrm{\Lambda }^{}{}_{\mu _1\mathrm{}\mu _7}{}^{\alpha \beta }=A_{\mu _1\mathrm{}\mu _7\sigma }^{\alpha \beta }\xi ^\sigma $$
(55)
provided the duality relation
$$F_{\mu _1\mathrm{}\mu _9}^{\alpha \beta }=iฯต_{\mu _1\mathrm{}\mu _9}{}_{}{}^{\sigma }[V_+^\alpha V_+^\beta P_\sigma ^{}V_{}^\alpha V_{}^\beta P_\sigma ]$$
(56)
holds, where $`F_{\mu _1\mathrm{}\mu _9}^{\alpha \beta }=9_{[\mu _1}A_{\mu _2\mathrm{}\mu _9]}^{\alpha \beta }+\mathrm{}`$, and the dots stand for terms involving lower rank forms. From the field strengths of the eight-forms, one can define the $`SU(1,1)`$ invariant quantity
$$G_{\mu _1\mathrm{}\mu _9}=ฯต_{\alpha \gamma }ฯต_{\beta \delta }V_+^\alpha V_+^\beta F_{\mu _1\mathrm{}\mu _9}^{\gamma \delta },$$
(57)
with $`U(1)`$ charge $`+2`$, and its complex conjugate
$$G_{\mu _1\mathrm{}\mu _9}^{}=ฯต_{\alpha \gamma }ฯต_{\beta \delta }V_{}^\alpha V_{}^\beta F_{\mu _1\mathrm{}\mu _9}^{\gamma \delta }.$$
(58)
In terms of these objects, the duality relation of eq. (56) becomes
$$G_{\mu _1\mathrm{}\mu _9}=iฯต_{\mu _1\mathrm{}\mu _9\sigma }P^\sigma ,G_{\mu _1\mathrm{}\mu _9}^{}=iฯต_{\mu _1\mathrm{}\mu _9\sigma }P^\sigma .$$
(59)
One can define a third nine-form,
$$\stackrel{~}{G}_{\mu _1\mathrm{}\mu _9}=ฯต_{\alpha \gamma }ฯต_{\beta \delta }V_+^\alpha V_{}^\beta F_{\mu _1\mathrm{}\mu _9}^{\gamma \delta },$$
(60)
with vanishing $`U(1)`$ charge, but the duality relation (56) implies that this nine-form vanishes identically , thus determining an $`SU(1,1)`$ invariant constraint. Therefore only two eight-forms are actually independent.
We now come to our choice for the field strengths, for which the most general general expression is
$`F_{\mu _1\mathrm{}\mu _9}^{\alpha \beta }=`$ $`9_{[\mu _1}A_{\mu _2\mathrm{}\mu _9]}^{\alpha \beta }+\alpha F_{[\mu _1\mathrm{}\mu _7}^{(\alpha }A_{\mu _8\mu _9]}^{\beta )}+\beta F_{[\mu _1\mathrm{}\mu _3}^{(\alpha }A_{\mu _4\mathrm{}\mu _9]}^{\beta )}+\gamma F_{[\mu _1\mathrm{}\mu _5}A_{\mu _6\mu _7}^{(\alpha }A_{\mu _8\mu _9]}^{\beta )}`$ (61)
$`+i\delta ฯต_{\gamma \delta }A_{[\mu _1\mu _2}^\gamma F_{\mu _3\mathrm{}\mu _5}^\delta A_{\mu _6\mu _7}^{(\alpha }A_{\mu _8\mu _9]}^{\beta )}+\xi A_{[\mu _1\mathrm{}\mu _4}F_{\mu _5\mathrm{}\mu _7}^{(\alpha }A_{\mu _8\mu _9]}^{\beta )}.`$
The freedom of redefining the eight-form, $`A_8A_8+A_6A_2+A_4A_2A_2`$, can be used to put to zero the coefficients $`\xi `$ and $`\delta `$ in (61). It turns out that defining the gauge transformation of the eight-forms as
$$\delta A_{\mu _1\mathrm{}\mu _8}^{\alpha \beta }=8_{[\mu _1}\mathrm{\Lambda }_{\mu _2\mathrm{}\mu _8]}^{\alpha \beta }+\frac{2}{9}\alpha F_{[\mu _1\mathrm{}\mu _7}^{(\alpha }\mathrm{\Lambda }_{\mu _8]}^{\beta )}+\frac{2}{3}\beta F_{[\mu _1\mathrm{}\mu _3}^{(\alpha }\mathrm{\Lambda }_{\mu _4\mathrm{}\mu _8]}^{\beta )},$$
(62)
the field strengths of eq. (61) are gauge invariant if the coefficient $`\gamma `$ vanishes as well, and if the coefficients $`\alpha `$ and $`\beta `$ are related by
$$\beta =7\alpha .$$
(63)
To summarize, we have obtained
$`F_{\mu _1\mathrm{}\mu _9}^{\alpha \beta }`$ $`=`$ $`9_{[\mu _1}A_{\mu _2\mathrm{}\mu _9]}^{\alpha \beta }+\alpha F_{[\mu _1\mathrm{}\mu _7}^{(\alpha }A_{\mu _8\mu _9]}^{\beta )}7\alpha F_{\mu _1\mathrm{}\mu _3}^{(\alpha }A_{\mu _4\mathrm{}\mu _9]}^{\beta )},`$ (64)
$`\delta A_{\mu _1\mathrm{}\mu _8}^{\alpha \beta }`$ $`=`$ $`8_{[\mu _1}\mathrm{\Lambda }_{\mu _2\mathrm{}\mu _8]}^{\alpha \beta }+\frac{2}{9}\alpha F_{[\mu _1\mathrm{}\mu _7}^{(\alpha }\mathrm{\Lambda }_{\mu _8]}^{\beta )}\frac{14}{3}\alpha F_{[\mu _1\mathrm{}\mu _3}^{(\alpha }\mathrm{\Lambda }_{\mu _4\mathrm{}\mu _8]}^{\beta )}.`$ (65)
We now consider the terms in the commutator coming from the last two lines of eq. (50), as well as the first two terms in eq. (52) and the part of the third containing lower rank forms. All these terms have to produce the gauge transformations of eq. (65) with the parameters given in eqs. (27) and (48), plus possibly a gauge transformation. The end result is that one produces the additional gauge transformation
$`\mathrm{\Lambda }_{\mu _1\mathrm{}\mu _7}^{{}_{}{}^{\prime \prime }\alpha \beta }`$ $`=`$ $`4ic\left[A_{[\mu _1\mathrm{}\mu _6}^{(\alpha }(V_+^{\beta )}\overline{ฯต}_2\gamma _{\mu _7]}ฯต_1+V_{}^{\beta )}\overline{ฯต}_{2C}\gamma _{\mu _7]}ฯต_{1C})\right]`$ (66)
$`+`$ $`12d\left[A_{[\mu _2\mu _3}^{(\alpha }(V_+^{\beta )}\overline{ฯต}_2\gamma _{\mu _4\mathrm{}\mu _7]}ฯต_1V_{}^{\beta )}\overline{ฯต}_{2C}\gamma _{\mu _4\mathrm{}\mu _7]}ฯต_{1C})\right],`$
and the algebra closes provided that the coefficients are fixed to be
$`c={\displaystyle \frac{21}{4}},d={\displaystyle \frac{7}{4}},e={\displaystyle \frac{105}{8}},`$
$`f=35,g=70,\alpha ={\displaystyle \frac{9}{4}}.`$ (67)
In conclusion, the supersymmetry transformation for the eight-forms is
$`\delta A_{\mu _1\mathrm{}\mu _8}^{\alpha \beta }`$ $`=`$ $`iV_+^\alpha V_+^\beta \overline{ฯต}\gamma _{\mu _1\mathrm{}\mu _8}\lambda _C+iV_{}^\alpha V_{}^\beta \overline{ฯต}_C\gamma _{\mu _1\mathrm{}\mu _8}\lambda `$ (68)
$`+`$ $`8V_+^{(\alpha }V_{}^{\beta )}\overline{ฯต}\gamma _{[\mu _1\mathrm{}\mu _7}\psi _{\mu _8]}8V_+^{(\alpha }V_{}^{\beta )}\overline{ฯต}_C\gamma _{[\mu _1\mathrm{}\mu _7}\psi _{\mu _8]C}`$
$`+`$ $`\frac{21}{4}A_{[\mu _1\mathrm{}\mu _6}^{(\alpha }\delta A_{\mu _7\mu _8]}^{\beta )}\frac{7}{4}A_{[\mu _1\mu _2}^{(\alpha }\delta A_{\mu _3\mathrm{}\mu _8]}^{\beta )}\frac{105i}{8}A_{[\mu _1\mu _2}^{(\alpha }A_{\mu _3\mu _4}^{\beta )}ฯต_{\gamma \delta }A_{\mu _5\mu _6}^\gamma \delta A_{\mu _7\mu _8]}^\delta `$
$``$ $`35A_{[\mu _1\mu _2}^{(\alpha }A_{\mu _3\mu _4}^{\beta )}\delta A_{\mu _5\mathrm{}\mu _8]}+70A_{[\mu _1\mathrm{}\mu _4}A_{\mu _5\mu _6}^{(\alpha }\delta A_{\mu _7\mu _8]}^{\beta )},`$
while the gauge invariance of the field strengths
$$F_{\mu _1\mathrm{}\mu _9}^{\alpha \beta }=9_{[\mu _1}A_{\mu _2\mathrm{}\mu _9]}^{\alpha \beta }+\frac{9}{4}F_{[\mu _1\mathrm{}\mu _7}^{(\alpha }A_{\mu _8\mu _9]}^{\beta )}\frac{63}{4}F_{[\mu _1\mathrm{}\mu _3}^{(\alpha }A_{\mu _4\mathrm{}\mu _9]}^{\beta )}$$
(69)
requires
$$\delta A_{\mu _1\mathrm{}\mu _8}^{\alpha \beta }=8_{[\mu _1}\mathrm{\Lambda }_{\mu _2\mathrm{}\mu _8]}^{\alpha \beta }+\frac{1}{2}F_{[\mu _1\mathrm{}\mu _7}^{(\alpha }\mathrm{\Lambda }_{\mu _8]}^{\beta )}\frac{21}{2}F_{[\mu _1\mathrm{}\mu _3}^{(\alpha }\mathrm{\Lambda }_{\mu _4\mathrm{}\mu _8]}^{\beta )}.$$
(70)
Finally, the seven-form gauge parameter that appears in the commutator is
$`\mathrm{\Lambda }_{\mu _1\mathrm{}\mu _7}^{\alpha \beta }`$ $`=`$ $`A_{\mu _1\mathrm{}\mu _7\sigma }^{\alpha \beta }\xi ^\sigma \frac{1}{2}S^{\alpha \beta }[\overline{ฯต}_2\gamma _{\mu _1\mathrm{}\mu _7}ฯต_1\overline{ฯต}_{2C}\gamma _{\mu _1\mathrm{}\mu _7}ฯต_{1C}]`$ (71)
$``$ $`\frac{21i}{8}A_{[\mu _1\mathrm{}\mu _6}^\alpha (V_+^\beta \overline{ฯต}_2\gamma _{\mu _7]}ฯต_1+V_{}^\beta \overline{ฯต}_{2C}\gamma _{\mu _7]}ฯต_{1C})`$
$``$ $`\frac{21}{8}A_{[\mu _1\mu _2}^\alpha (V_+^\beta \overline{ฯต}_2\gamma _{\mu _3\mathrm{}\mu _7]}ฯต_1V_{}^\beta \overline{ฯต}_{2C}\gamma _{\mu _3\mathrm{}\mu _7]}ฯต_{1C}),`$
as one obtains from eqs. (51), (55) and (66).
## 4 Ten-forms
The construction of ten-forms differs in an essential way from that of the six- and eight-forms: they do not have a field strength and therefore they cannot be dual to some other form within the IIB theory. They do not have propagating degrees of freedom, since the charge associated to them must vanish. Therefore there is no a priori limit on the number of ten-forms one could introduce. Also the $`SU(1,1)`$ representations cannot be guessed from the duality relations with lower rank forms. However, their supersymmetry transformations are well defined. We therefore proceed as before, determining the independent ten-forms from the requirement that the supersymmetry algebra must close. We want to determine the most general supersymmetry transformations for the ten-forms, compatible with $`U(1)`$ invariance, for a given $`SU(1,1)`$ representation. We first prove that both a doublet and a quadruplet of ten-forms are allowed, and then we discuss the claim that these are the only possible ten-forms that are compatible with all the symmetries of IIB supergravity.
### 4.1 The doublet of ten-forms
We want to determine the supersymmetry transformations of a doublet of ten-forms $`A_{\mu _1\mathrm{}\mu _{10}}^\alpha `$ satisfying the reality condition
$$(A^1)_{\mu _1\mathrm{}\mu _{10}}^{}=A_{\mu _1\mathrm{}\mu _{10}}^2.$$
(72)
As we have seen already in the previous sections, the supersymmetry transformation of any form consists of terms containing spinors, plus possibly terms containing lower-rank forms and their supersymmetry transformations. In the case of the ten-form doublet, $`U(1)`$ invariance requires that the most general fermionic part in the supersymmetry transformation of the ten-form doublet is
$`\delta A_{\mu _1\mathrm{}\mu _{10}}^\alpha `$ $`=`$ $`aV_{}^\alpha \overline{ฯต}\gamma _{\mu _1\mathrm{}\mu _{10}}\lambda +a^{}V_+^\alpha \overline{ฯต}_C\gamma _{\mu _1\mathrm{}\mu _{10}}\lambda _C`$ (73)
$`+`$ $`bV_{}^\alpha \overline{ฯต}_C\gamma _{[\mu _1\mathrm{}\mu _9}\psi _{\mu _{10}]}b^{}V_+^\alpha \overline{ฯต}\gamma _{[\mu _1\mathrm{}\mu _9}\psi _{\mu _{10}]C}.`$
The commutator of two such transformations contains the ten-form gauge transformation
$`\delta A_{\mu _1\mathrm{}\mu _{10}}^\alpha `$ $`=`$ $`10_{[\mu _1}\mathrm{\Lambda }_{\mu _2\mathrm{}\mu _{10}]}^\alpha `$ (74)
$`=`$ $`20i_{[\mu _1}\left(aV_+^\alpha \overline{ฯต}_2\gamma _{\mu _2\mathrm{}\gamma _{10}]}ฯต_{1C}+a^{}V_{}^\alpha \overline{ฯต}_{2C}\gamma _{\mu _2\mathrm{}\gamma _{10}]}ฯต_1\right),`$
provided that the coefficients $`a`$ and $`b`$ satisfy
$$b=20ia^{}.$$
(75)
Moreover, the additional terms in the commutator, containing the five-form $`F_5`$ and the complex three-form $`G_3`$, vanish if $`a`$ is chosen to be real.
In order to close the algebra, one also has to produce a general coordinate transformation with parameter $`\xi ^\mu `$ (27), but this exactly cancels with the gauge transformation of parameter<sup>6</sup><sup>6</sup>6For lower rank $`p`$-forms these transformations are obtained in the form $`\xi ^\rho F_{\rho \mu _1\mathrm{}\mu _p}`$, for $`p=D`$ the vanishing of the $`D+1`$-form $`F`$ corresponds to the cancellation of the two transformations. This result will be used again in the next subsection, when we will consider ten-forms in other representations of $`SU(1,1)`$.
$$\mathrm{\Lambda }_{\mu _1\mathrm{}\mu _9}^{}_{}{}^{}\alpha =A_{\mu _1\mathrm{}\mu _9\sigma }^\alpha \xi ^\sigma .$$
(76)
As a result, the algebra closes without adding any term containing lower-rank forms in the supersymmetry transformation of eq. (73). Correspondingly, this ten-form doublet is invariant with respect to the gauge transformations of the lower-rank forms. Without loss of generality, we can fix
$$a=1,$$
(77)
so that the resulting supersymmetry transformation for the ten-form doublet is
$`\delta A_{\mu _1\mathrm{}\mu _{10}}^\alpha `$ $`=`$ $`V_{}^\alpha \overline{ฯต}\gamma _{\mu _1\mathrm{}\mu _{10}}\lambda +V_+^\alpha \overline{ฯต}_C\gamma _{\mu _1\mathrm{}\mu _{10}}\lambda _C`$ (78)
$`+`$ $`20iV_{}^\alpha \overline{ฯต}_C\gamma _{[\mu _1\mathrm{}\mu _9}\psi _{\mu _{10}]}+20iV_+^\alpha \overline{ฯต}\gamma _{[\mu _1\mathrm{}\mu _9}\psi _{\mu _{10}]C}.`$
### 4.2 The quadruplet of ten-forms
We consider now a quadruplet of ten-forms $`A_{\mu _1\mathrm{}\mu _{10}}^{\alpha \beta \gamma }`$, completely symmetric in $`\alpha `$, $`\beta `$ and $`\gamma `$, and satisfying the reality condition
$$(A^{111})_{\mu _1\mathrm{}\mu _{10}}^{}=A_{\mu _1\mathrm{}\mu _{10}}^{222},(A^{112})_{\mu _1\mathrm{}\mu _{10}}^{}=A_{\mu _1\mathrm{}\mu _{10}}^{122}.$$
(79)
The most general supersymmetry transformation, compatible with the reality condition and with $`U(1)`$ invariance, and consisting of terms that only involve the spinors and terms containing the lower rank forms and their supersymmetry transformations, is
$`\delta A_{\mu _1\mathrm{}\mu _{10}}^{\alpha \beta \gamma }`$ $`=`$ $`aV_+^{(\alpha }V_+^\beta V_{}^{\gamma )}\overline{ฯต}_C\gamma _{\mu _1\mathrm{}\mu _{10}}\lambda _C+a^{}V_{}^{(\alpha }V_{}^\beta V_+^{\gamma )}\overline{ฯต}\gamma _{\mu _1\mathrm{}\mu _{10}}\lambda `$
$`+`$ $`bV_+^{(\alpha }V_+^\beta V_{}^{\gamma )}\overline{ฯต}\gamma _{[\mu _1\mathrm{}\mu _9}\psi _{\mu _{10}]C}b^{}V_{}^{(\alpha }V_{}^\beta V_+^{\gamma )}\overline{ฯต}_C\gamma _{[\mu _1\mathrm{}\mu _9}\psi _{\mu _{10}]}`$
$`+`$ $`cA_{[\mu _1\mathrm{}\mu _8}^{(\alpha \beta }\delta A_{\mu _9\mu _{10}]}^{\gamma )}+dA_{[\mu _1\mu _2}^{(\alpha }\delta A_{\mu _3\mathrm{}\mu _{10}]}^{\beta \gamma )}+eA_{[\mu _1\mathrm{}\mu _6}^{(\alpha }A_{\mu _7\mu _8}^\beta \delta A_{\mu _9\mu _{10}]}^{\gamma )}`$
$`+`$ $`fA_{[\mu _1\mu _2}^{(\alpha }A_{\mu _3\mu _4}^\beta \delta A_{\mu _5\mathrm{}\mu _{10}]}^{\gamma )}+gA_{[\mu _1\mathrm{}\mu _4}A_{\mu _5\mu _6}^{(\alpha }A_{\mu _7\mu _8}^\beta \delta A_{\mu _9\mu _{10}]}^{\gamma )}`$
$`+`$ $`hA_{[\mu _1\mu _2}^{(\alpha }A_{\mu _3\mu _4}^\beta A_{\mu _5\mu _6}^{\gamma )}\delta A_{\mu _7\mathrm{}\mu _{10}]}+ikA_{[\mu _1\mu _2}^{(\alpha }A_{\mu _3\mu _4}^\beta A_{\mu _5\mu _6}^{\gamma )}ฯต_{\delta \tau }A_{\mu _7\mu _8}^\delta \delta A_{\mu _9\mu _{10}]}^\tau .`$
We want to analyze the commutator of two such transformations.
We first consider the contribution coming from the fermionic terms, i.e., the first two lines of eq. (4.2). Those produce the ten-form gauge transformation
$`\delta A_{\mu _1\mathrm{}\mu _{10}}^{\alpha \beta \gamma }=10_{[\mu _1}\mathrm{\Lambda }_{\mu _2\mathrm{}\mu _{10}]}^{\alpha \beta \gamma }`$
$`=\frac{20}{3}i_{[\mu _1}(aV_+^{(\alpha }V_+^\beta V_{}^{\gamma )}\overline{ฯต}_2\gamma _{\mu _2\mathrm{}\mu _{10}]}ฯต_{1C}+a^{}V_{}^{(\alpha }V_{}^\beta V_+^{\gamma )}\overline{ฯต}_{2C}\gamma _{\mu _2\mathrm{}\mu _{10}]}ฯต_1)`$ (81)
together with the terms
$`\frac{20}{3}aF_{[\mu _1\mathrm{}\mu _9}^{(\alpha \beta }\left(V_+^{\gamma )}\overline{ฯต}_2\gamma _{\mu _{10}]}ฯต_{1C}+V_{}^{\gamma )}\overline{ฯต}_{2C}\gamma _{\mu _{10}]}ฯต_1\right)`$
$`20aiS^{(\alpha \beta }\left(\overline{ฯต}_2\gamma _{[\mu _1\mathrm{}\mu _7}ฯต_1\overline{ฯต}_{2C}\gamma _{[\mu _1\mathrm{}\mu _7}ฯต_{1C}\right)F_{\mu _8\mu _9\mu _{10}]}^{\gamma )},`$ (82)
provided that
$$\frac{20i}{3}a=b$$
(83)
and $`a`$ is chosen to be imaginary. Without loss of generality, we can fix
$$a=i$$
(84)
from now on. As in the case of the ten-form doublet of the previous subsection, a general coordinate transformation is automatically produced by means of a compensating gauge transformation of parameter
$$\mathrm{\Lambda }_{\mu _1\mathrm{}\mu _9}^{{}_{}{}^{}\alpha \beta \gamma }=A_{\mu _1\mathrm{}\mu _9\sigma }^{\alpha \beta \gamma }\xi ^\sigma .$$
(85)
We assume that the ten-form quadruplet transforms non-trivially with respect to the lower-rank form gauge transformations, and in particular we make the Ansatz
$$\delta A_{\mu _1\mathrm{}\mu _{10}}^{\alpha \beta \gamma }=\alpha F_{[\mu _1\mathrm{}\mu _9}^{(\alpha \beta }\mathrm{\Lambda }_{\mu _{10}]}^{\gamma )}+\beta F_{[\mu _1\mu _2\mu _3}^{(\alpha }\mathrm{\Lambda }_{\mu _4\mathrm{}\mu _{10}]}^{\beta \gamma )}.$$
(86)
We will comment on this choice at the end of this subsection. We now proceed exactly as in the previous cases, considering the terms in the commutator coming from the last three lines of eq. (4.2), as well as the two terms in eq. (82). Those have to generate the gauge transformations of eq. (86), possibly together with an additional ten-form gauge transformation. The final result is that the ten-form gauge transformation of parameter
$`\mathrm{\Lambda }_{\mu _1\mathrm{}\mu _9}^{{}_{}{}^{\prime \prime }\alpha \beta \gamma }`$ $`=`$ $`\frac{2i}{5}cA_{[\mu _1\mathrm{}\mu _8}^{(\alpha \beta }\left(V_+^{\gamma )}\overline{ฯต}_2\gamma _{\mu _9]}ฯต_{1C}+V_{}^{\gamma )}\overline{ฯต}_{2C}\gamma _{\mu _9]}ฯต_1\right)`$ (87)
$``$ $`\frac{2}{5}dA_{[\mu _1\mu _2}^{(\alpha }S^{\beta \gamma )}\left(\overline{ฯต}_2\gamma _{\mu _3\mathrm{}\mu _9]}ฯต_1\overline{ฯต}_{2C}\gamma _{\mu _3\mathrm{}\mu _9]}ฯต_{1C}\right)`$
is produced, while the coefficients are determined to be
$`\alpha ={\displaystyle \frac{2}{3}},\beta =32,c=12,`$
$`d=3,e={\displaystyle \frac{63}{4}},f={\displaystyle \frac{21}{4}},`$
$`g=210,h=105,k={\displaystyle \frac{315}{8}}.`$ (88)
Summarizing, the supersymmetry transformation of the ten-form quadruplet is
$`\delta A_{\mu _1\mathrm{}\mu _{10}}^{\alpha \beta \gamma }`$ $`=`$ $`iV_+^{(\alpha }V_+^\beta V_{}^{\gamma )}\overline{ฯต}_C\gamma _{\mu _1\mathrm{}\mu _{10}}\lambda _CiV_{}^{(\alpha }V_{}^\beta V_+^{\gamma )}\overline{ฯต}\gamma _{\mu _1\mathrm{}\mu _{10}}\lambda `$
$`+`$ $`\frac{20}{3}V_+^{(\alpha }V_+^\beta V_{}^{\gamma )}\overline{ฯต}\gamma _{[\mu _1\mathrm{}\mu _9}\psi _{\mu _{10}]C}\frac{20}{3}V_{}^{(\alpha }V_{}^\beta V_+^{\gamma )}\overline{ฯต}_C\gamma _{[\mu _1\mathrm{}\mu _9}\psi _{\mu _{10}]}`$
$``$ $`12A_{[\mu _1\mathrm{}\mu _8}^{(\alpha \beta }\delta A_{\mu _9\mu _{10}]}^{\gamma )}+3A_{[\mu _1\mu _2}^{(\alpha }\delta A_{\mu _3\mathrm{}\mu _{10}]}^{\beta \gamma )}\frac{63}{4}A_{[\mu _1\mathrm{}\mu _6}^{(\alpha }A_{\mu _7\mu _8}^\beta \delta A_{\mu _9\mu _{10}]}^{\gamma )}`$
$`+`$ $`\frac{21}{4}A_{[\mu _1\mu _2}^{(\alpha }A_{\mu _3\mu _4}^\beta \delta A_{\mu _5\mathrm{}\mu _{10}]}^{\gamma )}210A_{[\mu _1\mathrm{}\mu _4}A_{\mu _5\mu _6}^{(\alpha }A_{\mu _7\mu _8}^\beta \delta A_{\mu _9\mu _{10}]}^{\gamma )}`$
$`+`$ $`105A_{[\mu _1\mu _2}^{(\alpha }A_{\mu _3\mu _4}^\beta A_{\mu _5\mu _6}^{\gamma )}\delta A_{\mu _7\mathrm{}\mu _{10}]}+\frac{315i}{8}A_{[\mu _1\mu _2}^{(\alpha }A_{\mu _3\mu _4}^\beta A_{\mu _5\mu _6}^{\gamma )}ฯต_{\delta \tau }A_{\mu _7\mu _8}^\delta \delta A_{\mu _9\mu _{10}]}^\tau ,`$
while its gauge transformation is
$$\delta A_{\mu _1\mathrm{}\mu _{10}}^{\alpha \beta \gamma }=10_{[\mu _1}\mathrm{\Lambda }_{\mu _2\mathrm{}\mu _{10}]}^{\alpha \beta \gamma }\frac{2}{3}F_{[\mu _1\mathrm{}\mu _9}^{(\alpha \beta }\mathrm{\Lambda }_{\mu _{10}]}^{\gamma )}+32F_{[\mu _1\mu _2\mu _3}^{(\alpha }\mathrm{\Lambda }_{\mu _4\mathrm{}\mu _{10}]}^{\beta \gamma )}.$$
(90)
Finally, the ten-form gauge transformation parameter appearing in the supersymmetry algebra is
$`\mathrm{\Lambda }_{\mu _1\mathrm{}\mu _9}^{\alpha \beta \gamma }`$ $`=`$ $`A_{\mu _1\mathrm{}\mu _9\sigma }^{\alpha \beta \gamma }\xi ^\sigma \frac{2}{3}(V_+^{(\alpha }V_+^\beta V_{}^{\gamma )}\overline{ฯต}_2\gamma _{\mu _1\mathrm{}\mu _9}ฯต_{1C}V_{}^{(\alpha }V_{}^\beta V_+^{\gamma )}\overline{ฯต}_{2C}\gamma _{\mu _1\mathrm{}\mu _9}ฯต_1)`$ (91)
$`+`$ $`\frac{24i}{5}A_{[\mu _1\mathrm{}\mu _8}^{(\alpha \beta }\left(V_+^{\gamma )}\overline{ฯต}_2\gamma _{\mu _9]}ฯต_{1C}+V_{}^{\gamma )}\overline{ฯต}_{2C}\gamma _{\mu _9]}ฯต_1\right)`$
$``$ $`\frac{6}{5}A_{[\mu _1\mu _2}^{(\alpha }S^{\beta \gamma )}\left(\overline{ฯต}_2\gamma _{\mu _3\mathrm{}\mu _9]}ฯต_1\overline{ฯต}_{2C}\gamma _{\mu _3\mathrm{}\mu _9]}ฯต_{1C}\right),`$
as it results from eqs. (81), (85) and (87).
To conclude this subsection, we want to comment on the bosonic gauge transformation of eq. (90). Even though the supersymmetry algebra restricts us in our case to ten dimensions, it turns out that the bosonic gauge algebra closes for arbitrary dimension. In particular one can write down an eleven-form field strength that is gauge invariant with respect to a gauge transformation of the form (86):
$$F_{\mu _1\mathrm{}\mu _{11}}^{\alpha \beta \gamma }=11_{[\mu _1}A_{\mu _2\mathrm{}\mu _{11}]}^{\alpha \beta \gamma }+\frac{11}{2}\alpha A_{[\mu _1\mu _2}^{(\alpha }F_{\mu _3\mathrm{}\mu _{11}]}^{\beta \gamma )}+\frac{11}{8}\beta A_{[\mu _1\mathrm{}\mu _8}^{(\alpha \beta }F_{\mu _9\mu _{10}\mu _{11}]}^{\gamma )},$$
(92)
where the coefficients $`\alpha `$ and $`\beta `$ have to satisfy the constraint
$$\beta =48\alpha .$$
(93)
This relation is in agreement with the values of $`\alpha `$ and $`\beta `$ given in eq. (88) and obtained imposing supersymmetry. This suggests that the bosonic gauge algebra has an underlying structure that is independent of supersymmetry in ten dimensions<sup>7</sup><sup>7</sup>7This type of gauge algebra is also observed in the doubled fields approach, see ..
### 4.3 Other ten-forms?
We now want to show that no other ten-forms can be included in the supersymmetry algebra of IIB supergravity. In order to do this, we consider the most general Ansatz for the supersymmetry transformation of a ten-form in a generic representation of $`SU(1,1)`$. Without loss of generality, we can limit ourselves to ten-forms with vanishing $`U(1)`$-charge. The simplest such example is a singlet of $`SU(1,1)`$, for which the supersymmetry transformation necessarily is
$$\delta A_{\mu _1\mathrm{}\mu _{10}}=\overline{ฯต}\gamma _{[\mu _1\mathrm{}\mu _9}\psi _{\mu _{10}]}+\overline{ฯต}_C\gamma _{[\mu _1\mathrm{}\mu _9}\psi _{\mu _{10}]C}.$$
(94)
The commutator of two such transformations closes. This is not surprising since $`A_{(10)}`$ is the volume form,
$$A_{\mu _1\mathrm{}\mu _{10}}ฯต_{\mu _1\mathrm{}\mu _{10}}=e_{\mu _1}{}_{}{}^{a_1}\mathrm{}e_{\mu _{10}}{}_{}{}^{a_{10}}ฯต_{a_1\mathrm{}a_{10}}^{}.$$
(95)
This means that there are no independent ten-form singlets in the supersymmetry algebra of IIB.
One could ask whether additional ten-form doublets could result from objects of the form $`A_{\mu _1\mathrm{}\mu _{10}}^{\alpha _1\mathrm{}\alpha _{2n+1}}`$, when $`2n`$ $`SU(1,1)`$ indices are pairwise antisymmetrized. However, because of the constraint of eq. (3) these forms are the same as the one we obtained in section 4.1, and therefore there is only a single doublet of ten-forms in the theory. This argument can be iterated, so that for each object with an odd number of $`SU(1,1)`$ indices, only the components in the completely symmetric representation are independent of the ten-forms belonging to lower representations.
Therefore, given a ten-form with $`n`$ $`SU(1,1)`$ indices, one has to consider only the completely symmetric $`SU(1,1)`$ representation. Let us consider the case $`n=2`$ first. The most general Ansatz for the fermionic terms is
$$\delta A_{(10)}^{\alpha \beta }=a_1V_+^{(\alpha }V_{}^{\beta )}\overline{ฯต}\gamma _{(9)}\psi +a_2V_+^{(\alpha }V_{}^{\beta )}\overline{ฯต}_C\gamma _{(9)}\psi _C+b_1V_+^\alpha V_+^\beta \overline{ฯต}\gamma _{(10)}\lambda _C+b_2V_{}^\alpha V_{}^\beta \overline{ฯต}_C\gamma _{(10)}\lambda .$$
(96)
As in the case of the singlet, one can close the algebra on this Ansatz, but again it is not an independent field. It turns out to be the variation of a composite field:
$$\delta \left(\frac{1}{2}S^{\alpha \beta }ฯต_{(10)}\right)=\delta \left(V_+^{(\alpha }V_{}^{\beta )}ฯต_{(10)}\right)=V_+^{(\alpha }V_{}^{\beta )}\delta ฯต_{(10)}+\delta V_+^{(\alpha }V_{}^{\beta )}ฯต_{(10)}+V_+^{(\alpha }\delta V_{}^{\beta )}ฯต_{(10)}.$$
(97)
This generalises to ten-forms with $`n=2m`$ $`SU(1,1)`$-indices, for which we can also close the algebra, but end up with the variation of the composite field
$$S^{(\alpha _1\beta _1}\mathrm{}S^{\alpha _m\beta _m)}ฯต_{(10)}.$$
(98)
The case of $`n`$ odd is different, since the requirement of vanishing $`U(1)`$ charge does not allow one to write down a volume form. In this case the Ansatz for the fermionic part of the supersymmetry transformation is (we set here $`n=2m+1`$)
$`\delta A_{\mu _1\mathrm{}\mu _{10}}^{\alpha _1\mathrm{}\alpha _{2m+1}}`$ $`=`$ $`aV_+^{(\alpha _1}\mathrm{}V_+^{\alpha _{m+1}}V_{}^{\alpha _{m+2}}\mathrm{}V_{}^{\alpha _{2m+1})}\overline{ฯต}_C\gamma _{\mu _1\mathrm{}\mu _{10}}\lambda _C`$ (99)
$`+`$ $`a^{}V_{}^{(\alpha _1}\mathrm{}V_{}^{\alpha _{m+1}}V_+^{\alpha _{m+2}}\mathrm{}V_+^{\alpha _{2m+1})}\overline{ฯต}\gamma _{\mu _1\mathrm{}\mu _{10}}\lambda `$
$`+`$ $`bV_+^{(\alpha _1}\mathrm{}V_+^{\alpha _{m+1}}V_{}^{\alpha _{m+2}}\mathrm{}V_{}^{\alpha _{2m+1})}\overline{ฯต}\gamma _{[\mu _1\mathrm{}\mu _9}\psi _{\mu _{10}]C}`$
$``$ $`b^{}V_{}^{(\alpha _1}\mathrm{}V_{}^{\alpha _{m+1}}V_+^{\alpha _{m+2}}\mathrm{}V_+^{\alpha _{2m+1})}\overline{ฯต}_C\gamma _{[\mu _1\mathrm{}\mu _9}\psi _{\mu _{10}]}.`$
It can be shown that only for the case $`m=0`$, i.e., the doublet that we already considered, the commutator of two such transformations closes producing just a ten-form gauge transformation and a general coordinate transformation. As we have seen already for the quadruplet ($`m=1`$), extra terms are generated that need to combine with additional terms in eq. (99), containing lower-rank forms and their supersymmetry transformations, to produce bosonic gauge transformations. An explicit analysis shows that these terms can only be written for the quadruplet. Higher $`SU(1,1)`$ representations require introducing additional contributions from the scalars in these bosonic terms, and the supersymmetry commutator produces derivatives of these scalars. These contributions can not be identified with any parameter that appears in the supersymmetry algebra. This suggests that only a doublet and a quadruplet can be consistently included in the supersymmetry algebra of IIB.
## 5 The complete IIB transformation rules and algebra
This section collects our results for the $`SU(1,1)`$-democratic version of $`D=10`$ IIB supergravity. We present the supersymmetry transformation rules, the transformation rules of the $`p`$-forms under bosonic gauge transformations, the definition of gauge invariant curvatures, and finally the results for the commutator of two supersymmetry transformations. Of course all the transformations and definitions are interdependent. All results have been derived only up to the quadratic order in the fermions.
The supersymmetry transformation rules in Einstein frame, in the notation of , are:
$`\delta e_\mu ^a`$ $`=`$ $`i\overline{ฯต}\gamma ^a\psi _\mu +i\overline{ฯต}_C\gamma ^a\psi _{\mu C},`$ (100)
$`\delta \psi _\mu `$ $`=`$ $`D_\mu ฯต+\frac{i}{480}F_{\mu \nu _1\mathrm{}\nu _4}\gamma ^{\nu _1\mathrm{}\nu _4}ฯต+\frac{1}{96}G^{\nu \rho \sigma }\gamma _{\mu \nu \rho \sigma }ฯต_C\frac{3}{32}G_{\mu \nu \rho }\gamma ^{\nu \rho }ฯต_C,`$ (101)
$`\delta A_{\mu \nu }^\alpha `$ $`=`$ $`V_{}^\alpha \overline{ฯต}\gamma _{\mu \nu }\lambda +V_+^\alpha \overline{ฯต}_C\gamma _{\mu \nu }\lambda _C+4iV_{}^\alpha \overline{ฯต}_C\gamma _{[\mu }\psi _{\nu ]}+4iV_+^\alpha \overline{ฯต}\gamma _{[\mu }\psi _{\nu ]C},`$ (102)
$`\delta A_{\mu \nu \rho \sigma }`$ $`=`$ $`\overline{ฯต}\gamma _{[\mu \nu \rho }\psi _{\sigma ]}\overline{ฯต}_C\gamma _{[\mu \nu \rho }\psi _{\sigma ]C}\frac{3i}{8}ฯต_{\alpha \beta }A_{[\mu \nu }^\alpha \delta A_{\rho \sigma ]}^\beta ,`$ (103)
$`\delta \lambda `$ $`=`$ $`iP_\mu \gamma ^\mu ฯต_C\frac{i}{24}G_{\mu \nu \rho }\gamma ^{\mu \nu \rho }ฯต,`$ (104)
$`\delta V_+^\alpha `$ $`=`$ $`V_{}^\alpha \overline{ฯต}_C\lambda ,`$ (105)
$`\delta V_{}^\alpha `$ $`=`$ $`V_+^\alpha \overline{ฯต}\lambda _C,`$ (106)
$`\delta A_{\mu _1\mathrm{}\mu _6}^\alpha `$ $`=`$ $`iV_{}^\alpha \overline{ฯต}\gamma _{\mu _1\mathrm{}\mu _6}\lambda iV_+^\alpha \overline{ฯต}_C\gamma _{\mu _1\mathrm{}\mu _6}\lambda _C`$ (107)
$`+12\left(V_{}^\alpha \overline{ฯต}_C\gamma _{[\mu _1\mathrm{}\mu _5}\psi _{\mu _6]}V_+^\alpha \overline{ฯต}\gamma _{[\mu _1\mathrm{}\mu _5}\psi _{C\mu _6]}\right)`$
$`+40A_{[\mu _1\mathrm{}\mu _4}\delta A_{\mu _5\mu _6]}^\alpha 20\delta A_{[\mu _1\mathrm{}\mu _4}A_{\mu _5\mu _6]}^\alpha `$
$`\frac{15i}{2}A_{[\mu _1\mu _2}^\alpha ฯต_{\beta \gamma }A_{\mu _3\mu _4}^\beta \delta A_{\mu _5\mu _6]}^\gamma ,`$
$`\delta A_{\mu _1\mathrm{}\mu _8}^{\alpha \beta }`$ $`=`$ $`+iV_{}^{(\alpha }V_{}^{\beta )}\overline{ฯต}_C\gamma _{\mu _1\mathrm{}\mu _8}\lambda iV_+^{(\alpha }V_+^{\beta )}\overline{ฯต}\gamma _{\mu _1\mathrm{}\mu _8}\lambda _C`$ (108)
$`+8V_+^{(\alpha }V_{}^{\beta )}(\overline{ฯต}\gamma _{[\mu _1\mathrm{}\mu _7}\psi _{\mu _8]}\overline{ฯต}_C\gamma _{[\mu _1\mathrm{}\mu _7}\psi _{C\mu _8]})`$
$`+\frac{21}{4}A_{[\mu _1\mathrm{}\mu _6}^{(\alpha }\delta A_{\mu _7\mu _8]}^{\beta )}\frac{7}{4}A_{[\mu _1\mu _2}^{(\alpha }\delta A_{\mu _3\mathrm{}\mu _8]}^{\beta )}`$
$`35A_{[\mu _1\mu _2}^{(\alpha }A_{\mu _3\mu _4}^{\beta )}\delta A_{\mu _5\mathrm{}\mu _8]}+70A_{[\mu _1\mathrm{}\mu _4}A_{\mu _5\mu _6}^{(\alpha }\delta A_{\mu _7\mu _8]}^{\beta )}`$
$`\frac{105i}{8}A_{[\mu _1\mu _2}^{(\alpha }A_{\mu _3\mu _4}^{\beta )}ฯต_{\gamma \delta }A_{\mu _5\mu _6}^\gamma \delta A_{\mu _7\mu _8]}^\delta ,`$
$`\delta A_{\mu _1\mathrm{}\mu _{10}}^\alpha `$ $`=`$ $`V_{}^\alpha \overline{ฯต}\gamma _{\mu _1\mathrm{}\mu _{10}}\lambda +V_+^\alpha \overline{ฯต}_C\gamma _{\mu _1\mathrm{}\mu _{10}}\lambda _C`$ (109)
$`+20i\left(V_+^\alpha \overline{ฯต}\gamma _{[\mu _1\mathrm{}\mu _9}\psi _{C\mu _{10}]}+V_{}^\alpha \overline{ฯต}_C\gamma _{[\mu _1\mathrm{}\mu _9}\psi _{\mu _{10}]}\right),`$
$`\delta A_{\mu _1\mathrm{}\mu _{10}}^{\alpha \beta \gamma }`$ $`=`$ $`iV_+^{(\alpha }V_+^\beta V_{}^{\gamma )}\overline{ฯต}_C\gamma _{\mu _1\mathrm{}\mu _{10}}\lambda _CiV_{}^{(\alpha }V_{}^\beta V_+^{\gamma )}\overline{ฯต}\gamma _{\mu _1\mathrm{}\mu _{10}}\lambda `$ (110)
$`+\frac{20}{3}(V_+^{(\alpha }V_+^\beta V_{}^{\gamma )}\overline{ฯต}\gamma _{[\mu _1\mathrm{}\mu _9}\psi _{C\mu _{10}]}V_{}^{(\alpha }V_{}^\beta V_+^{\gamma )}\overline{ฯต}_C\gamma _{[\mu _1\mathrm{}\mu _9}\psi _{\mu _{10}]})`$
$`12A_{[\mu _1\mathrm{}\mu _8}^{(\alpha \beta }\delta A_{\mu _9\mu _{10}]}^{\gamma )}+3A_{[\mu _1\mu _2}^{(\alpha }\delta A_{\mu _3\mathrm{}\mu _{10}]}^{\beta \gamma )}`$
$`\frac{63}{4}A_{[\mu _1\mathrm{}\mu _6}^{(\alpha }A_{\mu _7\mu _8}^\beta \delta A_{\mu _9\mu _{10}]}^{\gamma )}+\frac{21}{4}A_{[\mu _1\mu _2}^{(\alpha }A_{\mu _3\mu _4}^\beta \delta A_{\mu _5\mathrm{}\mu _{10}]}^{\gamma )}`$
$`\mathrm{\hspace{0.17em}210}A_{[\mu _1\mathrm{}\mu _4}A_{\mu _5\mu _6}^{(\alpha }A_{\mu _7\mu _8}^\beta \delta A_{\mu _9\mu _{10}]}^{\gamma )}+105A_{[\mu _1\mu _2}^{(\alpha }A_{\mu _3\mu _4}^\beta A_{\mu _5\mu _6}^{\gamma )}\delta A_{\mu _7\mathrm{}\mu _{10}]}`$
$`+\frac{315i}{8}A_{[\mu _1\mu _2}^{(\alpha }A_{\mu _3\mu _4}^\beta A_{\mu _5\mu _6}^{\gamma )}ฯต_{\delta \tau }A_{\mu _7\mu _8}^\delta \delta A_{\mu _9\mu _{10}]}^\tau .`$
For the bosonic gauge-transformations we find:
$`\delta A_{\mu _1\mu _2}^\alpha `$ $`=`$ $`2_{[\mu _1}\mathrm{\Lambda }_{\mu _2]}^\alpha ,`$ (111)
$`\delta A_{\mu _1\mathrm{}\mu _4}`$ $`=`$ $`4_{[\mu _1}\mathrm{\Lambda }_{\mu _2\mu _3\mu _4]}\frac{i}{4}ฯต_{\gamma \delta }\mathrm{\Lambda }_{[\mu _1}^\gamma F_{\mu _2\mu _3\mu _4]}^\delta ,`$ (112)
$`\delta A_{\mu _1\mathrm{}\mu _6}^\alpha `$ $`=`$ $`6_{[\mu _1}\mathrm{\Lambda }_{\mu _2\mathrm{}\mu _6]}^\alpha 8\mathrm{\Lambda }_{[\mu _1}^\alpha F_{\mu _2\mathrm{}\mu _6]}\frac{160}{3}F_{[\mu _1\mu _2\mu _3}^\alpha \mathrm{\Lambda }_{\mu _4\mu _5\mu _6]},`$ (113)
$`\delta A_{\mu _1\mathrm{}\mu _8}^{\alpha \beta }`$ $`=`$ $`8_{[\mu _1}\mathrm{\Lambda }_{\mu _2\mathrm{}\mu _8]}^{(\alpha \beta )}+\frac{1}{2}F_{[\mu _1\mathrm{}\mu _7}^{(\alpha }\mathrm{\Lambda }_{\mu _8]}^{\beta )}\frac{21}{2}F_{[\mu _1\mu _2\mu _3}^{(\alpha }\mathrm{\Lambda }_{\mu _4\mathrm{}\mu _8]}^{\beta )},`$ (114)
$`\delta A_{\mu _1\mathrm{}\mu _{10}}^\alpha `$ $`=`$ $`10_{[\mu _1}\mathrm{\Lambda }_{\mu _2\mathrm{}\mu _{10}]}^\alpha ,`$ (115)
$`\delta A_{\mu _1\mathrm{}\mu _{10}}^{\alpha \beta \gamma }`$ $`=`$ $`10_{[\mu _1}\mathrm{\Lambda }_{\mu _2\mathrm{}\mu _{10}]}^{(\alpha \beta \gamma )}\frac{2}{3}F_{[\mu _1\mathrm{}\mu _9}^{(\alpha \beta }\mathrm{\Lambda }_{\mu _{10}]}^{\gamma )}+32F_{[\mu _1\mu _2\mu _3}^{(\alpha }\mathrm{\Lambda }_{\mu _4\mathrm{}\mu _{10}]}^{\beta \gamma )}.`$ (116)
For the $`p`$-form fields we define field-strengths invariant under the bosonic gauge transformations
$`F_{\mu _1\mu _2\mu _3}^\alpha `$ $`=`$ $`3_{[\mu _1}A_{\mu _2\mu _3]}^\alpha ,`$ (117)
$`F_{\mu _1\mathrm{}\mu _5}`$ $`=`$ $`5_{[\mu _1}A_{\mu _2\mathrm{}\mu _5]}+\frac{5i}{8}ฯต_{\alpha \beta }A_{[\mu _1\mu _2}^\alpha F_{\mu _3\mu _4\mu _5]}^\beta ,`$ (118)
$`F_{\mu _1\mathrm{}\mu _7}^\alpha `$ $`=`$ $`7_{[\mu _1}A_{\mu _2\mathrm{}\mu _7]}^\alpha +28A_{[\mu _1\mu _2}^\alpha F_{\mu _3\mathrm{}\mu _7]}\frac{280}{3}F_{[\mu _1\mu _2\mu _3}^\alpha A_{\mu _4\mathrm{}\mu _7]},`$ (119)
$`F_{\mu _1\mathrm{}\mu _9}^{\alpha \beta }`$ $`=`$ $`9_{[\mu _1}A_{\mu _2\mathrm{}\mu _9]}^{\alpha \beta }+\frac{9}{4}F_{[\mu _1\mathrm{}\mu _7}^{(\alpha }A_{\mu _8\mu _9]}^{\beta )}\frac{63}{4}F_{[\mu _1\mu _2\mu _3}^{(\alpha }A_{\mu _4\mathrm{}\mu _9]}^{\beta )},`$ (120)
$`F_{\mu _1\mathrm{}\mu _{11}}^\alpha `$ $`=`$ $`11_{[\mu _1}A_{\mu _2\mathrm{}\mu _{11}]}^\alpha =0,`$ (121)
$`F_{\mu _1\mathrm{}\mu _{11}}^{\alpha \beta \gamma }`$ $`=`$ $`11(_{[\mu _1}A_{\mu _2\mathrm{}\mu _{11}]}^{\alpha \beta \gamma }\frac{1}{3}F_{[\mu _1\mathrm{}\mu _9}^{(\alpha \beta }A_{\mu _{10}\mu _{11}]}^{\gamma )}+4F_{[\mu _1\mu _2\mu _3}^{(\alpha }A_{\mu _4\mathrm{}\mu _{11}[}^{\beta \gamma )})=0.`$ (122)
The duality relations between these field-strengths are:
$`F_{\mu _1\mathrm{}\mu _7}^\alpha `$ $`=`$ $`\frac{i}{3!}ฯต_{\mu _1\mathrm{}\mu _7\mu \nu \rho }S^{\alpha \beta }ฯต_{\beta \gamma }F^{\gamma ;\mu \nu \rho },`$ (123)
$`F_{\mu _1\mathrm{}\mu _9}^{\alpha \beta }`$ $`=`$ $`iฯต_{\mu _1\mathrm{}\mu _9}{}_{}{}^{\rho }[V_+^\alpha V_+^\beta P_\rho ^{}V_{}^\alpha V_{}^\beta P_\rho ].`$ (124)
The commutator of two supersymmetry transformations, $`[\delta (ฯต_1),\delta (ฯต_2)]`$ must close on symmetry transformations of the IIB multiplet. In fact, as we saw in previous sections, this is the way the results of this paper have been obtained. We find the following contributions to $`[\delta (ฯต_1),\delta (ฯต_2)]`$:
$`\xi ^\mu `$ $`=`$ $`i\overline{ฯต}_2\gamma ^\mu ฯต_1+i\overline{ฯต}_{2C}\gamma ^\mu ฯต_{1C},`$ (125)
$`\mathrm{\Lambda }_\mu ^\alpha `$ $`=`$ $`A_{\mu \nu }^\alpha \xi ^\nu 2i[V_+^\alpha \overline{ฯต}_2\gamma _\mu ฯต_{1C}+V_{}^\alpha \overline{ฯต}_{2C}\gamma _\mu ฯต_1],`$ (126)
$`\mathrm{\Lambda }_{\mu _1\mu _2\mu _3}`$ $`=`$ $`A_{\mu _1\mu _2\mu _3\nu }\xi ^\nu +\frac{1}{4}[\overline{ฯต}_2\gamma _{\mu _1\mu _2\mu _3}ฯต_1\overline{ฯต}_{2C}\gamma _{\mu _1\mu _2\mu _3}ฯต_{1C}],`$ (127)
$`\mathrm{\Lambda }_{\mu _1\mathrm{}\mu _5}^\alpha `$ $`=`$ $`A_{\mu _1\mathrm{}\mu _5\rho }^\alpha \xi ^\rho 2V_{}^\alpha \overline{ฯต}_{2C}\gamma _{\mu _1\mathrm{}\mu _5}ฯต_1+2V_+^\alpha \overline{ฯต}_2\gamma _{\mu _1\mathrm{}\mu _5}ฯต_{1C}`$ (128)
$`+\frac{40}{3}A_{[\mu _1\mathrm{}\mu _4}\mathrm{\Lambda }_{\mu _5]}^\alpha \frac{40}{3}\mathrm{\Lambda }_{[\mu _1\mu _2\mu _3}A_{\mu _4\mu _5]}^\alpha ,`$
$`\mathrm{\Lambda }_{\mu _1\mathrm{}\mu _7}^{\alpha \beta }`$ $`=`$ $`A_{\mu _1\mathrm{}\mu _7\nu }^{(\alpha \beta )}\xi ^\nu 2V_+^{(\alpha }V_{}^{\beta )}\left(\overline{ฯต}_2\gamma _{\mu _1\mathrm{}\mu _7}ฯต_1\overline{ฯต}_{2C}\gamma _{\mu _1\mathrm{}\mu _7}ฯต_{1C}\right)`$ (129)
$`+\frac{21}{16}A_{[\mu _1\mathrm{}\mu _6}^{(\alpha }\mathrm{\Lambda }_{\mu _7]}^{\beta )}\frac{21}{16}A_{[\mu _1\mu _2}^{(\alpha }\mathrm{\Lambda }_{\mu _3\mathrm{}\mu _7]}^{\beta )},`$
$`\mathrm{\Lambda }_{\mu _1\mathrm{}\mu _9}^\alpha `$ $`=`$ $`2i\left(V_+^\alpha \overline{ฯต}_2\gamma _{\mu _1\mathrm{}\mu _9}ฯต_{1C}+V_{}^\alpha \overline{ฯต}_{2C}\gamma _{\mu _1\mathrm{}\mu _9}ฯต_1\right)`$ (130)
$`\mathrm{\Lambda }_{\mu _1\mathrm{}\mu _9}^{\alpha \beta \gamma }`$ $`=`$ $`A_{\mu _1\mathrm{}\mu _9\nu }^{(\alpha \beta \gamma )}\xi ^\nu \frac{2}{3}\left(V_+^{(\alpha }V_+^\beta V_{}^{\gamma )}\overline{ฯต}_2\gamma _{\mu _1\mathrm{}\mu _9}ฯต_{1C}V_{}^{(\alpha }V_{}^\beta V_+^{\gamma )}\overline{ฯต}_{2C}\gamma _{\mu _1\mathrm{}\mu _9}ฯต_1\right)`$ (131)
$`+\frac{24i}{5}A_{[\mu _1\mathrm{}\mu _8}^{(\alpha \beta }(V_+^{\gamma )}\overline{ฯต}_2\gamma _{\mu _9]}ฯต_{1C}+V_{}^{\gamma )}\overline{ฯต}_{2C}\gamma _{\mu _9]}ฯต_1)`$
$`\frac{6}{5}A_{[\mu _1\mu _2}^{(\alpha }S^{\beta \gamma )}(\overline{ฯต}_2\gamma _{\mu _3\mathrm{}\mu _9]}ฯต_1\overline{ฯต}_{2C}\gamma _{\mu _3\mathrm{}\mu _9]}ฯต_{1C}).`$
This concludes the summary of our main results. In the next section we will present the IIB supergravity multiplet in a real formulation in both Einstein frame and string frame.
## 6 $`U(1)`$ gauge fixing and string frame
The results we derived so far were in Einstein frame. To go to string frame we will first choose a $`U(1)`$ gauge, so that the dependence on the dilaton becomes explicit. Our choice is
$$V_{}^1V_+^2=V_{}^1.$$
(132)
To preserve this condition the supersymmetry transformations have to be modified by a field dependent $`U(1)`$ gauge transformation:
$$\delta ^{}(ฯต)=\delta (ฯต)+\delta _{U(1)}\left(\frac{i}{2V_{}^1}(V_{}^2\overline{ฯต}_C\lambda V_+^1\overline{ฯต}\lambda _C)\right).$$
(133)
This modification is only visible on the scalars since on the fermions it gives rise to terms cubic in fermionic variables. The $`SU(1,1)`$ transformations are also modified: the condition (132) is preserved under a combination of an $`SU(1,1)`$ and a $`U(1)`$ transformation. On a field $`\chi `$ of $`U(1)`$-charge $`q`$ the required $`U(1)`$ transformation is
$$\chi e^{iq\theta }\chi ,\mathrm{with}e^{2i\theta }=\frac{\alpha +\beta z}{\overline{\alpha }+\overline{\beta }\overline{z}},$$
(134)
where the coordinate $`z`$ is defined in (17). This is of course visible on all fermions.
To make the dilaton and axion explicit we set
$$V_{}^1=V_+^2=\frac{1}{\sqrt{1z\overline{z}}},V_{}^2=(V_+^1)^{}=\frac{z}{\sqrt{1z\overline{z}}}.$$
(135)
Using (21) and (25) we find (we now drop the prime on the redefined supersymmetry transformation)
$$\delta \tau =2ie^\varphi e^{2i\mathrm{\Lambda }}\overline{ฯต}\lambda _C,$$
(136)
where
$$e^{2i\mathrm{\Lambda }}=\frac{1i\tau }{1+i\overline{\tau }}.$$
(137)
Useful variables are
$$P_\mu =\frac{i}{2}e^\varphi e^{2i\mathrm{\Lambda }}_\mu \tau ,Q_\mu =\frac{1}{4}e^\varphi \left(\frac{1i\overline{\tau }}{1i\tau }_\mu \tau +\frac{1+i\tau }{1+i\overline{\tau }}_\mu \overline{\tau }\right).$$
(138)
It is convenient to get rid of the factors of $`e^{2i\mathrm{\Lambda }}`$ in the supersymmetry transformation rules . To do this we redefine the fermions by phase factors according to their $`U(1)`$ weights:
$$\lambda ^{}=e^{3i\mathrm{\Lambda }/2}\lambda ,\psi _\mu ^{}=e^{i\mathrm{\Lambda }/2}\psi _\mu ,ฯต^{}=e^{i\mathrm{\Lambda }/2}ฯต.$$
(139)
In the transformation rules the scalars $`V_\pm ^\alpha `$ will now occur everywhere in the combination $`V_\pm ^\alpha e^{\pm i\mathrm{\Lambda }}`$, which are:
$`V_{}^1e^{i\mathrm{\Lambda }}`$ $`=`$ $`{\displaystyle \frac{1}{2}}e^{\varphi /2}(1i\tau ),`$
$`V_+^1e^{i\mathrm{\Lambda }}`$ $`=`$ $`{\displaystyle \frac{1}{2}}e^{\varphi /2}(1i\overline{\tau }),`$
$`V_{}^2e^{i\mathrm{\Lambda }}`$ $`=`$ $`{\displaystyle \frac{1}{2}}e^{\varphi /2}(1+i\tau ),`$ (140)
$`V_+^2e^{i\mathrm{\Lambda }}`$ $`=`$ $`{\displaystyle \frac{1}{2}}e^{\varphi /2}(1+i\overline{\tau }).`$
Note that interchanging $`V^1V^2`$ corresponds to $`\tau \tau `$, $`V_+V_{}`$ to $`\tau \overline{\tau }`$. The transformation rules for the IIB supergravity multiplet of now become :
$`\delta e_\mu ^a`$ $`=`$ $`i(\overline{ฯต}\gamma ^a\psi _\mu )+\mathrm{h}.\mathrm{c}.`$ (141)
$`\delta \psi _\mu `$ $`=`$ $`D_\mu ฯต\frac{i}{4}e^\varphi ฯต_\mu \mathrm{}+\frac{i}{480}F_{\mu _1\mathrm{}\mu _5}\gamma ^{\mu _1\mathrm{}\mu _5}\gamma _\mu ฯต`$ (142)
$`\frac{1}{192}e^{\varphi /2}\left(\gamma _\mu \gamma ^{\nu \rho \sigma }+2\gamma ^{\nu \rho \sigma }\gamma _\mu \right)ฯต_C(F^1F^2+i\overline{\tau }(F^1+F^2))_{\nu \rho \sigma },`$
$`\delta A_{\mu \nu }^1`$ $`=`$ $`2ie^{\varphi /2}\{(1i\overline{\tau })(\overline{ฯต}\gamma _{[\mu }\psi _{C\nu ]}\frac{i}{4}\overline{ฯต}_C\gamma _{\mu \nu }\lambda _C)`$ (143)
$`+(1i\tau )(\overline{ฯต}_C\gamma _{[\mu }\psi _{\nu ]}\frac{i}{4}ฯต\gamma _{\mu \nu ]}\lambda )\},`$
$`\delta A_{\mu \nu }^2`$ $`=`$ $`2ie^{\varphi /2}\{(1+i\overline{\tau })(\overline{ฯต}\gamma _{[\mu }\psi _{C\nu ]}\frac{i}{4}\overline{ฯต}_C\gamma _{\mu \nu }\lambda _C)`$ (144)
$`+(1+i\tau )(\overline{ฯต}_C\gamma _{[\mu }\psi _{\nu ]}\frac{i}{4}ฯต\gamma _{\mu \nu ]}\lambda )\},`$
$`\delta A_{\mu \nu \rho \sigma }`$ $`=`$ $`\overline{ฯต}\gamma _{[\mu \nu \rho }\psi _{\sigma ]}+\mathrm{h}.\mathrm{c}.\frac{3i}{8}ฯต_{\alpha \beta }A_{[\mu \nu }^\alpha \delta A_{\rho \sigma ]}^\beta ,`$ (145)
$`\delta \lambda `$ $`=`$ $`\frac{1}{2}e^\varphi \gamma ^\mu ฯต_C_\mu \overline{\tau }\frac{i}{48}e^{\varphi /2}\gamma ^{\nu \rho \sigma }ฯต(F^1F^2+i\overline{\tau }(F^1+F^2))_{\nu \rho \sigma },`$ (146)
$`\delta \mathrm{}`$ $`=`$ $`ie^\varphi (\overline{ฯต}_C\lambda \overline{ฯต}\lambda _C),`$ (147)
$`\delta \varphi `$ $`=`$ $`\overline{ฯต}_C\lambda +\overline{ฯต}\lambda _C.`$ (148)
For the higher-rank form fields we present only the transformations to the fermions, because the contributions containing explicit gauge fields are unchanged by the gauge fixing and redefinitions and can be read off from (107-110). We find:
$`\delta A^1_{(6)}`$ $`=`$ $`6e^{\varphi /2}\{(1i\tau )(\overline{ฯต}_C\gamma _{(5)}\psi +\frac{i}{12}\overline{ฯต}\gamma _{(6)}\lambda )`$ (149)
$`(1i\overline{\tau })(\overline{ฯต}\gamma _{(5)}\psi +\frac{i}{12}\overline{ฯต}_C\gamma _{(6)}\lambda _C)\}+\mathrm{},`$
$`\delta A^2_{(6)}`$ $`=`$ $`6e^{\varphi /2}\{(1+i\tau )(\overline{ฯต}_C\gamma _{(5)}\psi +\frac{i}{12}\overline{ฯต}\gamma _{(6)}\lambda )`$ (150)
$`(1+i\overline{\tau })(\overline{ฯต}\gamma _{(5)}\psi _C+\frac{i}{12}\overline{ฯต}_C\gamma _{(6)}\lambda _C)\}+\mathrm{},`$
$`\delta A^{11}_{(8)}`$ $`=`$ $`2e^\varphi \{(1i\tau )(1i\overline{\tau })(\overline{ฯต}\gamma _{(7)}\psi \overline{ฯต}_C\gamma _{(7)}\psi _C)`$ (151)
$`+\frac{i}{8}((1i\tau )^2\overline{ฯต}_C\gamma _{(8)}\lambda (1i\overline{\tau })^2\overline{ฯต}\gamma _{(8)}\lambda _C)\}+\mathrm{},`$
$`\delta A^{22}_{(8)}`$ $`=`$ $`2e^\varphi \{(1+i\tau )(1+i\overline{\tau })(\overline{ฯต}\gamma _{(7)}\psi \overline{ฯต}_C\gamma _{(7)}\psi _C)`$ (152)
$`+\frac{i}{8}((1+i\tau )^2\overline{ฯต}_C\gamma _{(8)}\lambda (1+i\overline{\tau })^2\overline{ฯต}\gamma _{(8)}\lambda _C)\}+\mathrm{},`$
$`\delta A^{12}_{(8)}`$ $`=`$ $`e^\varphi \{((1i\tau )(1+i\overline{\tau })+(1+i\tau )(1i\overline{\tau }))(\overline{ฯต}\gamma _{(7)}\psi \overline{ฯต}_C\gamma _{(7)}\psi _C)`$ (153)
$`+\frac{i}{4}((1i\tau )(1+i\tau )\overline{ฯต}_C\gamma _{(8)}\lambda (1i\overline{\tau })(1+i\overline{\tau })\overline{ฯต}\gamma _{(8)}\lambda _C)\}+\mathrm{},`$
$`\delta A^1_{(10)}`$ $`=`$ $`10ie^{\varphi /2}\{(1i\overline{\tau })(\overline{ฯต}\gamma _{(9)}\psi _C\frac{i}{20}\overline{ฯต}_C\gamma _{(10)}\lambda _C)`$ (154)
$`+(1i\tau )(\overline{ฯต}_C\gamma _{(9)}\psi \frac{i}{20}\overline{ฯต}\gamma _{(10)}\lambda )\},`$
$`\delta A^2_{(10)}`$ $`=`$ $`10ie^{\varphi /2}\{(1+i\overline{\tau })(\overline{ฯต}\gamma _{(9)}\psi _C\frac{i}{20}\overline{ฯต}_C\gamma _{(10)}\lambda _C)`$ (155)
$`+(1+i\tau )(\overline{ฯต}_C\gamma _{(9)}\psi \frac{i}{20}\overline{ฯต}\gamma _{(10)}\lambda )\},`$
$`\delta A_{(10)}^{111}`$ $`=`$ $`\frac{5}{6}e^{3\varphi /2}(1i\tau )(1i\overline{\tau })\{(1i\overline{\tau })(\overline{ฯต}\gamma _{(9)}\psi _C+\frac{3i}{20}\overline{ฯต}_C\gamma _{(10)}\lambda _C)`$ (156)
$`(1i\tau )(\overline{ฯต}_C\gamma _{(9)}\psi +\frac{3i}{20}\overline{ฯต}\gamma _{(10)}\lambda )\}+\mathrm{},`$
$`\delta A_{(10)}^{222}`$ $`=`$ $`\frac{5}{6}e^{3\varphi /2}(1+i\tau )(1+i\overline{\tau })\{(1+i\overline{\tau })(\overline{ฯต}\gamma _{(9)}\psi _C+\frac{3i}{20}\overline{ฯต}_C\gamma _{(10)}\lambda _C)`$ (157)
$`(1+i\tau )(\overline{ฯต}_C\gamma _{(9)}\psi +\frac{3i}{20}\overline{ฯต}\gamma _{(10)}\lambda )\}+\mathrm{},`$
$`\delta A_{(10)}^{112}`$ $`=`$ $`\frac{5}{18}e^{3\varphi /2}\{((1i\overline{\tau })^2(1+i\tau )+2(1i\tau )(1i\overline{\tau })(1+i\overline{\tau }))\times `$ (158)
$`\times (\overline{ฯต}\gamma _{(9)}\psi _C+\frac{3i}{20}\overline{ฯต}_C\gamma _{(10)}\lambda _C)`$
$`((1i\tau )^2(1+i\overline{\tau })+2(1i\tau )(1+i\tau )(1i\overline{\tau }))\times `$
$`\times (\overline{ฯต}_C\gamma _{(9)}\psi +\frac{3i}{20}\overline{ฯต}\gamma _{(10)}\lambda )\}+\mathrm{},`$
$`\delta A_{(10)}^{221}`$ $`=`$ $`\frac{5}{18}e^{3\varphi /2}\{((1+i\overline{\tau })^2(1i\tau )+2(1+i\tau )(1+i\overline{\tau })(1i\overline{\tau }))\times `$ (159)
$`\times (\overline{ฯต}\gamma _{(9)}\psi _C+\frac{3i}{20}\overline{ฯต}_C\gamma _{(10)}\lambda _C)`$
$`((1+i\tau )^2(1i\overline{\tau })+2(1+i\tau )(1i\tau )(1+i\overline{\tau }))\times `$
$`\times (\overline{ฯต}_C\gamma _{(9)}\psi +\frac{3i}{20}\overline{ฯต}\gamma _{(10)}\lambda )\}+\mathrm{},`$
Here the dots stand for the gauge field terms given in (107-110). So in formulas (141) to (159) we have collected the complete set of Einstein frame supersymmetry transformations in the real formulation.
Let us now review the transformations under $`SU(1,1)`$ and $`SL(2,)`$ transformations. Consider an $`SU(1,1)`$ transformation
$$U=\left(\begin{array}{cc}\alpha & \beta \\ \overline{\beta }& \overline{\alpha }\end{array}\right),\alpha \overline{\alpha }\beta \overline{\beta }=1.$$
(160)
The field $`\tau `$ transforms under the corresponding $`SL(2,)`$ transformation as
$`\tau `$ $``$ $`{\displaystyle \frac{a\tau +b}{c\tau +d}},\delta \tau {\displaystyle \frac{\delta \tau }{(c\tau +d)^2}},adbc=1,`$
$`a`$ $`=`$ $`\mathrm{Re}(\alpha \beta ),d=\mathrm{Re}(\alpha +\beta ),b=\mathrm{Im}(\alpha +\beta ),c=\mathrm{Im}(\alpha \beta ).`$ (161)
The redefinition (139) modifies the behavior under $`SU(1,1)`$ transformations. The compensating $`U(1)`$ transformation on a field $`\chi `$ of charge $`q`$ (134) is now changed to
$$\chi e^{iq\xi }\chi ,\mathrm{with}e^{2i\xi }=\frac{c\tau +d}{c\overline{\tau }+d}.$$
(162)
For the dilaton one finds
$$e^\varphi e^\varphi (c\tau +d)(c\overline{\tau }+d).$$
(163)
One easily verifies that, e.g., the supersymmetry variation of $`\tau `$
$$\delta \tau =2ie^\varphi \overline{ฯต}\lambda _C$$
(164)
is consistent with these transformations. The bosonic fields with vanishing $`U(1)`$ charge still transform in the standard way under $`SU(1,1)`$.
We will now bring some order into the collection of higher-rank forms (149) to (159) by considering certain linear combinations of these. We choose the linear combinations of the $`n`$-forms such that for a given $`n`$, each combination has a unique power of $`\tau `$ in the fermionic terms of the supersymmetry variation. This is motivated by the fact that the RR-forms come with a prefactor of $`e^\varphi `$ in the standard string frame basis, which is proportional to $`\tau \overline{\tau }`$. Thus we make the following definitions:
$`\stackrel{~}{C}_{(2)}`$ $`=`$ $`A_{(2)}^1A_{(2)}^2,\stackrel{~}{B}_{(2)}=A_{(2)}^1+A_{(2)}^2,`$ (165)
$`\stackrel{~}{C}_{(4)}`$ $`=`$ $`A_{(4)},`$ (166)
$`\stackrel{~}{C}_{(6)}`$ $`=`$ $`A_{(6)}^1+A_{(6)}^2,\stackrel{~}{B}_{(6)}=A_{(6)}^1A_{(6)}^2,`$ (167)
$`\stackrel{~}{C}_{(8)}`$ $`=`$ $`A_{(8)}^{11}+A_{(8)}^{22}+2A_{(8)}^{12},\stackrel{~}{B}_{(8)}=A_{(8)}^{11}+A_{(8)}^{22}2A_{(8)}^{12},`$ (168)
$`\stackrel{~}{D}_{(8)}`$ $`=`$ $`A_{(8)}^{11}A_{(8)}^{22},`$ (169)
$`\stackrel{~}{๐}_{(10)}`$ $`=`$ $`A_{(10)}^1+A_{(10)}^2,\stackrel{~}{}_{(10)}=A_{(10)}^1A_{(10)}^2,`$ (170)
$`\stackrel{~}{C}_{(10)}`$ $`=`$ $`A_{(10)}^{111}+A_{(10)}^{222}+3\left(A_{(10)}^{112}+A_{(10)}^{221}\right),`$ (171)
$`\stackrel{~}{B}_{(10)}`$ $`=`$ $`A_{(10)}^{111}A_{(10)}^{222}3\left(A_{(10)}^{112}A_{(10)}^{221}\right),`$ (172)
$`\stackrel{~}{E}_{(10)}`$ $`=`$ $`A_{(10)}^{111}+A_{(10)}^{222}\left(A_{(10)}^{112}+A_{(10)}^{221}\right),`$ (173)
$`\stackrel{~}{D}_{(10)}`$ $`=`$ $`A_{(10)}^{111}A_{(10)}^{222}+\left(A_{(10)}^{112}A_{(10)}^{221}\right).`$ (174)
A nice property, and partial justification why we refer to some of these linear combinations as $`\stackrel{~}{C}_{(n)}`$ ( RR fields) and $`\stackrel{~}{B}_{(n)}`$ ( NS-NS fields) is the way these fields transform into each other under $`S`$-duality. The discrete $`S`$-duality transformation $`\tau 1/\tau `$ corresponds to an $`SL(2,)`$-transformation with $`a=d=0`$, $`b=c=1`$. The behaviour of the form-fields under $`S`$-duality is
$`\stackrel{~}{C}_{(2)}`$ $``$ $`i\stackrel{~}{B}_{(2)},\stackrel{~}{B}_{(2)}i\stackrel{~}{C}_{(2)},`$
$`\stackrel{~}{C}_{(4)}`$ $``$ $`\stackrel{~}{C}_{(4)},`$
$`\stackrel{~}{C}_{(6)}`$ $``$ $`i\stackrel{~}{B}_{(6)},\stackrel{~}{B}_{(6)}i\stackrel{~}{C}_{(6)},`$
$`\stackrel{~}{C}_{(8)}`$ $``$ $`\stackrel{~}{B}_{(8)},\stackrel{~}{B}_{(8)}\stackrel{~}{C}_{(8)},\stackrel{~}{D}_{(8)}\stackrel{~}{D}_{(8)},`$
$`\stackrel{~}{๐}_{(10)}`$ $``$ $`i\stackrel{~}{}_{(10)},\stackrel{~}{}_{(10)}i\stackrel{~}{๐}_{(10)},`$
$`\stackrel{~}{C}_{(10)}`$ $``$ $`i\stackrel{~}{B}_{(10)},\stackrel{~}{B}_{(10)}i\stackrel{~}{C}_{(10)},`$
$`\stackrel{~}{D}_{(10)}`$ $``$ $`i\stackrel{~}{E}_{(10)},\stackrel{~}{E}_{(10)}i\stackrel{~}{D}_{(10)}.`$ (175)
We see that applying $`S`$-duality twice gives +1 on $`\tau `$ and on the four- and eight-forms, but $`1`$ on the two-, six- and ten-forms. That this is indeed right, and that the $`S`$-duality transformation is not its own inverse can be seen easily from translating back to the $`SU(1,1)`$ notation via (161), in which the $`S`$-duality transformation matrix is given by
$$U=\left(\begin{array}{cc}i& 0\\ 0& i\end{array}\right)$$
(176)
so that $`U^2`$ gives a minus on forms with an odd number of SU(1,1)-indices.
Now we are ready to transform to string frame. The basic transformation is $`e_{(E)\mu }{}_{}{}^{a}=e^{\varphi /4}e_{(S)\mu }^a`$. We choose to write the variation of the zehnbein in standard form, which requires a modification of supersymmetry with a $`\lambda `$-dependent local Lorentz transformation (which we see only on the zehnbein), and a redefinition:
$`ฯต^{}`$ $`=`$ $`e^{\varphi /8}ฯต,`$ (177)
$`\lambda ^{}`$ $`=`$ $`e^{\varphi /8}\lambda ,`$ (178)
$`\psi _\mu ^{}`$ $`=`$ $`e^{\varphi /8}\psi _\mu \frac{i}{4}\gamma _\mu ^{}\lambda _C^{},`$ (179)
$`\gamma _\mu ^{}`$ $`=`$ $`e^{\varphi /4}\gamma _\mu .`$ (180)
Again we start with the basic supergravity multiplet and then discuss the high-rank forms. The transformation rules are simplified by writing the complex fermions as real doublets, i.e. $`ฯต(ฯต_1,ฯต_2)`$, where $`ฯต_i`$ are real Majorana-Weyl fermions. This gives rise to the appearance of Pauli matrices $`\sigma _0=1,\sigma _1,i\sigma _2,\sigma _3`$ in the contractions between such doublets, generically:
$`\overline{ฯต}_C\gamma \chi +\overline{ฯต}\gamma \chi _C2\overline{ฯต}\sigma _3\gamma \chi ,\overline{ฯต}_C\gamma \chi _C+\overline{ฯต}\gamma \chi 2\overline{ฯต}\gamma \chi ,`$ (181)
$`\overline{ฯต}_C\gamma \chi \overline{ฯต}\gamma \chi _C2i\overline{ฯต}\sigma _1\gamma \chi ,\overline{ฯต}_C\gamma \chi _C\overline{ฯต}\gamma \chi 2i\overline{ฯต}(i\sigma _2)\gamma \chi .`$ (182)
In addition we redefine $`\lambda \lambda _C`$, or, equivalently, in the real notation
$$\lambda \sigma _3\lambda .$$
(183)
We drop all primes in the string frame transformation rules:
$`\delta e_\mu ^a`$ $`=`$ $`2i\overline{ฯต}\gamma ^a\psi _\mu `$ (184)
$`\delta \psi _\mu `$ $`=`$ $`D_\mu ฯต+\frac{1}{8}e^\varphi \gamma ^\nu _\nu \mathrm{}\gamma _\mu (i\sigma _2)ฯต\frac{1}{16}\gamma ^{\nu \rho }\sigma _3ฯตF_{+\mu \nu \rho }`$ (185)
$`+\frac{i}{96}e^\varphi \gamma ^{\nu \rho \sigma }\gamma _\mu \sigma _1ฯต(F_{}+i\mathrm{}F_+)_{\nu \rho \sigma }`$
$`\frac{1}{480}e^\varphi \gamma ^{\mu _1\mathrm{}\mu _5}\gamma _\mu (i\sigma _2)ฯตF_{\mu _1\mathrm{}\mu _5},`$
$`\delta \stackrel{~}{B}_{\mu \nu }`$ $`=`$ $`8i\overline{ฯต}\sigma _3\gamma _{[\mu }\psi _{\nu ]},`$ (186)
$`\delta \stackrel{~}{C}_{\mu \nu }`$ $`=`$ $`8e^\varphi \overline{ฯต}\sigma _1\gamma _{[\mu }(\psi _{\nu ]}+\frac{i}{2}\gamma _{\nu ]}\lambda )i\mathrm{}\delta \stackrel{~}{B}_{\mu \nu },`$ (187)
$`\delta \stackrel{~}{C}_{\mu \nu \rho \sigma }`$ $`=`$ $`2ie^\varphi \overline{ฯต}(i\sigma _2)\gamma _{[\mu \nu \rho }(\psi _{\sigma ]}+\frac{i}{4}\gamma _{\sigma ]}\lambda )`$ (188)
$`\frac{3i}{16}\{\stackrel{~}{C}\delta \stackrel{~}{B}\stackrel{~}{B}\delta \stackrel{~}{C}\}_{\mu \nu \rho \sigma },`$
$`\delta \lambda `$ $`=`$ $`\frac{i}{2}\gamma ^\mu _\mu \varphi ฯต\frac{i}{48}\gamma ^{\nu \rho \sigma }\sigma _3ฯตF_{+\nu \rho \sigma }\frac{i}{2}e^\varphi \gamma ^\mu _\mu \mathrm{}(i\sigma _2)ฯต`$ (189)
$`+\frac{1}{48}e^\varphi \gamma ^{\nu \rho \sigma }\sigma _1ฯต(F_{}+i\mathrm{}F_+)_{\nu \rho \sigma },`$
$`\delta \mathrm{}`$ $`=`$ $`2e^\varphi \overline{ฯต}(i\sigma _2)\lambda ,`$ (190)
$`\delta \varphi `$ $`=`$ $`2\overline{ฯต}\lambda `$ (191)
where we have defined
$$F_+=F^1+F^2,F_{}=F^1F^2.$$
(192)
For the higher form fields we find:
$`\delta \stackrel{~}{C}_{(6)}`$ $`=`$ $`24ie^\varphi \overline{ฯต}\sigma _1\gamma _{(5)}(\psi +\frac{i}{6}\gamma _{(1)}\lambda )+40\stackrel{~}{C}_{(4)}\delta \stackrel{~}{B}_{(2)}`$ (193)
$`20\delta \stackrel{~}{C}_{(4)}\stackrel{~}{B}_{(2)}\frac{15i}{4}\stackrel{~}{B}_{(2)}\left(\stackrel{~}{C}_{(2)}\delta \stackrel{~}{B}_{(2)}\stackrel{~}{B}_{(2)}\delta \stackrel{~}{C}_{(2)}\right),`$
$`\delta \stackrel{~}{B}_{(6)}`$ $`=`$ $`24e^\varphi \left\{\mathrm{}\overline{ฯต}\sigma _1\gamma _{(5)}(\psi +\frac{i}{6}\gamma _{(1)}\lambda )+e^\varphi \overline{ฯต}\sigma _3\gamma _{(5)}(\psi +\frac{i}{3}\gamma _{(1)}\lambda )\right\}`$ (194)
$`+40\stackrel{~}{C}_{(4)}\delta \stackrel{~}{C}_{(2)}20\delta \stackrel{~}{C}_{(4)}\stackrel{~}{C}_{(2)}\frac{15i}{4}\stackrel{~}{C}_{(2)}\left(\stackrel{~}{C}_{(2)}\delta \stackrel{~}{B}_{(2)}\stackrel{~}{B}_{(2)}\delta \stackrel{~}{C}_{(2)}\right),`$
$`\delta \stackrel{~}{C}_{(8)}`$ $`=`$ $`16ie^\varphi \overline{ฯต}(i\sigma _2)\gamma _{(7)}(\psi +\frac{i}{8}\gamma _{(1)}\lambda )+\frac{21}{4}\stackrel{~}{C}_{(6)}\delta \stackrel{~}{B}_{(2)}\frac{7}{4}\stackrel{~}{B}_{(2)}\delta \stackrel{~}{C}_{(6)}`$ (195)
$`35\stackrel{~}{B}_{(2)}\stackrel{~}{B}_{(2)}\delta \stackrel{~}{C}_{(4)}+70\stackrel{~}{C}_{(4)}\stackrel{~}{B}_{(2)}\delta \stackrel{~}{B}_{(2)]}`$
$`\frac{105i}{16}\stackrel{~}{B}_{(2)}\stackrel{~}{B}_{(2)}\left(\stackrel{~}{C}_{(2)}\delta \stackrel{~}{B}_{(2)}\stackrel{~}{B}_{(2)}\delta \stackrel{~}{C}_{(2)}\right),`$
$`\delta \stackrel{~}{B}_{(8)}`$ $`=`$ $`16ie^\varphi \{\mathrm{}^2\overline{ฯต}(i\sigma _2)\gamma _{(7)}(\psi +\frac{i}{8}\gamma _{(1)}\lambda )`$ (196)
$`+\frac{i}{4}\mathrm{}e^\varphi \overline{ฯต}\gamma _{(8)}\lambda +e^{2\varphi }\overline{ฯต}(i\sigma _2)\gamma _{(7)}(\psi +\frac{3i}{8}\gamma _{(1)}\lambda )\}`$
$`+\frac{21}{4}\stackrel{~}{B}_{(6)}\delta \stackrel{~}{C}_{(2)}\frac{7}{4}\stackrel{~}{C}_{(2)}\delta \stackrel{~}{B}_{(6)}35\stackrel{~}{C}_{(2)}\stackrel{~}{C}_{(2)}\delta \stackrel{~}{C}_{(4)}`$
$`+70\stackrel{~}{C}_{(4)}\stackrel{~}{C}_{(2)}\delta \stackrel{~}{C}_{(2)}\frac{105i}{16}\stackrel{~}{C}_{(2)}\stackrel{~}{C}_{(2)}\left(\stackrel{~}{C}_{(2)}\delta \stackrel{~}{B}_{(2)}\stackrel{~}{B}_{(2)}\delta \stackrel{~}{C}_{(2)}\right),`$
$`\delta \stackrel{~}{D}_{(8)}`$ $`=`$ $`16\mathrm{}e^\varphi \overline{ฯต}(i\sigma _2)\gamma _{(7)}(\psi +\frac{i}{8}\gamma _{(1)}\lambda )+2ie^{2\varphi }\overline{ฯต}\gamma _{(8)}\lambda `$ (197)
$`+\frac{21}{8}\{\stackrel{~}{C}_{(6)}\delta \stackrel{~}{C}_{(2)}+\stackrel{~}{B}_{(6)}\delta \stackrel{~}{B}_{(2)}\}`$
$`\frac{7}{8}\{\stackrel{~}{B}_{(2)}\delta \stackrel{~}{B}_{(6)}+\stackrel{~}{C}_{(2)}\delta \stackrel{~}{C}_{(6)}\}35\stackrel{~}{B}_{(2)}\stackrel{~}{C}_{(2)}\delta \stackrel{~}{C}_{(4)}`$
$`+35\stackrel{~}{C}_{(4)}\{\stackrel{~}{B}_{(2)}\delta \stackrel{~}{C}_{(2)}+\stackrel{~}{C}_{(2)}\delta \stackrel{~}{B}_{(2)}\}`$
$`\frac{105i}{16}\stackrel{~}{B}_{(2)}\stackrel{~}{C}_{(2)}\left(\stackrel{~}{C}_{(2)}\delta \stackrel{~}{B}_{(2)}\stackrel{~}{B}_{(2)}\delta \stackrel{~}{C}_{(2)}\right),`$
$`\delta \stackrel{~}{๐}_{(10)}`$ $`=`$ $`40ie^{2\varphi }\overline{ฯต}\sigma _3\gamma _{(9)}(\psi +\frac{i}{5}\gamma _{(1)}\lambda ),`$ (198)
$`\delta \stackrel{~}{}_{(10)}`$ $`=`$ $`40e^{2\varphi }\left\{\mathrm{}\overline{ฯต}\sigma _3\gamma _{(9)}(\psi +\frac{i}{5}\gamma _{(1)}\lambda )e^\varphi \overline{ฯต}\sigma _1\gamma _{(9)}(\psi +\frac{3i}{10}\gamma _{(1)}\lambda )\right\},`$ (199)
$`\delta \stackrel{~}{C}_{(10)}`$ $`=`$ $`\frac{40}{3}ie^\varphi \overline{ฯต}\sigma _1\gamma _{(9)}\left(\psi +\frac{i}{10}\gamma _{(1)}\lambda \right)`$ (201)
$`12\stackrel{~}{C}_{(8)}\delta \stackrel{~}{B}_{(2)}+3\delta \stackrel{~}{C}_{(8)}\stackrel{~}{B}_{(2)}\frac{63}{4}\stackrel{~}{C}_{(6)}\stackrel{~}{B}_{(2)}\delta \stackrel{~}{B}_{(2)}+\frac{21}{4}\delta \stackrel{~}{C}_{(6)}\stackrel{~}{B}_{(2)}\stackrel{~}{B}_{(2)}`$
$`210\stackrel{~}{C}_{(4)}\stackrel{~}{B}_{(2)}\stackrel{~}{B}_{(2)}\delta \stackrel{~}{B}_{(2)}+105\delta \stackrel{~}{C}_{(4)}\stackrel{~}{B}_{(2)}\stackrel{~}{B}_{(2)}\stackrel{~}{B}_{(2)}`$
$`+\frac{315}{16}i\stackrel{~}{B}_{(2)}\stackrel{~}{B}_{(2)}\stackrel{~}{B}_{(2)}\left(\stackrel{~}{C}_{(2)}\delta \stackrel{~}{B}_{(2)}\stackrel{~}{B}_{(2)}\delta \stackrel{~}{C}_{(2)}\right),`$
$`\delta \stackrel{~}{B}_{(10)}`$ $`=`$ $`+\frac{40}{3}e^{4\varphi }\left(1+e^{2\varphi }\mathrm{}^2\right)\overline{ฯต}\sigma _3\gamma _{(9)}\left(\psi +\frac{2i}{5}\gamma _{(1)}\lambda \right)`$ (202)
$`+\frac{40}{3}\mathrm{}e^\varphi \left(\mathrm{}^2+e^{2\varphi }\right)\overline{ฯต}\sigma _1\gamma _{(9)}\left(\psi +\frac{i}{10}\gamma _{(1)}\lambda \right)`$
$`12\stackrel{~}{B}_{(8)}\delta \stackrel{~}{C}_{(2)}+3\delta \stackrel{~}{B}_{(8)}\stackrel{~}{C}_{(2)}\frac{63}{4}\stackrel{~}{B}_{(6)}\stackrel{~}{C}_{(2)}\delta \stackrel{~}{C}_{(2)}+\frac{21}{4}\delta \stackrel{~}{B}_{(6)}\stackrel{~}{C}_{(2)}\stackrel{~}{C}_{(2)}`$
$`210\stackrel{~}{C}_{(4)}\stackrel{~}{C}_{(2)}\stackrel{~}{C}_{(2)}\delta \stackrel{~}{C}_{(2)}+105\delta \stackrel{~}{C}_{(4)}\stackrel{~}{C}_{(2)}\stackrel{~}{C}_{(2)}\stackrel{~}{C}_{(2)}`$
$`+\frac{315}{16}i\stackrel{~}{C}_{(2)}\stackrel{~}{C}_{(2)}\stackrel{~}{C}_{(2)}\left(\stackrel{~}{C}_{(2)}\delta \stackrel{~}{B}_{(2)}\stackrel{~}{B}_{(2)}\delta \stackrel{~}{C}_{(2)}\right),`$
$`\delta \stackrel{~}{D}_{(10)}`$ $`=`$ $`\frac{40}{9}e^{2\varphi }\overline{ฯต}\sigma _3\gamma _{(9)}\left(\psi +\frac{2i}{5}\gamma _{(1)}\lambda \right)\frac{40}{3}\mathrm{}e^\varphi \overline{ฯต}\sigma _1\gamma _{(9)}\left(\psi +\frac{i}{10}\gamma _{(1)}\lambda \right)`$
$`+2\stackrel{~}{B}_{(2)}\delta \stackrel{~}{D}_{(8)}+\stackrel{~}{C}_{(2)}\delta \stackrel{~}{C}_{(8)}8\stackrel{~}{D}_{(8)}\delta \stackrel{~}{B}_{(2)}4\stackrel{~}{C}_{(8)}\delta \stackrel{~}{C}_{(2)}`$
$`+\frac{7}{4}\left(\stackrel{~}{B}_{(2)}^2\delta \stackrel{~}{B}_{(6)}+2\stackrel{~}{B}_{(2)}\stackrel{~}{C}_{(2)}\delta \stackrel{~}{C}_{(6)}\right)`$
$`\frac{21}{4}\left(\stackrel{~}{B}_{(2)}\stackrel{~}{C}_{(6)}\delta \stackrel{~}{C}_{(2)}+\stackrel{~}{C}_{(2)}\stackrel{~}{C}_{(6)}\delta \stackrel{~}{B}_{(2)}+\stackrel{~}{B}_{(2)}\stackrel{~}{B}_{(6)}\delta \stackrel{~}{B}_{(2)}\right)`$
$`70\left(\stackrel{~}{B}_{(2)}^2\stackrel{~}{C}_{(4)}\delta \stackrel{~}{C}_{(2)}+2\stackrel{~}{B}_{(2)}\stackrel{~}{C}_{(2)}\stackrel{~}{C}_{(4)}\delta \stackrel{~}{B}_{(2)}\frac{3}{2}\stackrel{~}{B}_{(2)}^2\stackrel{~}{C}_{(2)}\delta \stackrel{~}{C}_{(4)}\right)`$
$`+\frac{315}{16}i\left(\stackrel{~}{B}_{(2)}^2\stackrel{~}{C}_{(2)}^2\delta \stackrel{~}{B}_{(2)}\stackrel{~}{B}_{(2)}^3\stackrel{~}{C}_{(2)}\delta \stackrel{~}{C}_{(2)}\right),`$
$`\delta \stackrel{~}{E}_{(10)}`$ $`=`$ $`+\frac{40}{3}ie^{3\varphi }\left(\frac{1}{3}+\mathrm{}^2e^{2\varphi }\right)\overline{ฯต}\sigma _1\gamma _{(9)}\left(\psi +\frac{i}{10}\gamma _{(1)}\lambda \right)`$ (203)
$`+\frac{80}{9}i\mathrm{}e^{2\varphi }\overline{ฯต}\sigma _3\gamma _{(9)}\left(\psi +\frac{2i}{5}\gamma _{(1)}\lambda \right)`$
$`+2\stackrel{~}{C}_{(2)}\delta \stackrel{~}{D}_{(8)}+\stackrel{~}{B}_{(2)}\delta \stackrel{~}{B}_{(8)}8\stackrel{~}{D}_{(8)}\delta \stackrel{~}{C}_{(2)}4\stackrel{~}{B}_{(8)}\delta \stackrel{~}{B}_{(2)}`$
$`+\frac{7}{4}\left(\stackrel{~}{C}_{(2)}^2\delta \stackrel{~}{C}_{(6)}+2\stackrel{~}{B}_{(2)}\stackrel{~}{C}_{(2)}\delta \stackrel{~}{B}_{(6)}\right)`$
$`\frac{21}{4}\left(\stackrel{~}{B}_{(2)}\stackrel{~}{B}_{(6)}\delta \stackrel{~}{C}_{(2)}+\stackrel{~}{C}_{(2)}\stackrel{~}{C}_{(6)}\delta \stackrel{~}{C}_{(2)}+\stackrel{~}{C}_{(2)}\stackrel{~}{B}_{(6)}\delta \stackrel{~}{B}_{(2)}\right)`$
$`70\left(\stackrel{~}{C}_{(2)}^2\stackrel{~}{C}_{(4)}\delta \stackrel{~}{B}_{(2)}+2\stackrel{~}{B}_{(2)}\stackrel{~}{C}_{(2)}\stackrel{~}{C}_{(4)}\delta \stackrel{~}{C}_{(2)}\frac{3}{2}\stackrel{~}{B}_{(2)}\stackrel{~}{C}_{(2)}^2\delta \stackrel{~}{C}_{(4)}\right)`$
$`+\frac{315}{16}i\left(\stackrel{~}{B}_{(2)}\stackrel{~}{C}_{(2)}^3\delta \stackrel{~}{B}_{(2)}\stackrel{~}{B}_{(2)}^2\stackrel{~}{C}_{(2)}^2\delta \stackrel{~}{C}_{(2)}\right).`$
We will now introduce the standard RR and NS-NS fields, and extend this to the higher rank forms. For this we define:
$`B_{(2)}`$ $`=`$ $`\frac{1}{2}\stackrel{~}{B}_{(2)},`$ (204)
$`C_{(0)}`$ $`=`$ $`\frac{1}{2}\mathrm{},`$ (205)
$`C_{(2)}`$ $`=`$ $`\frac{i}{4}\stackrel{~}{C}_{(2)},`$ (206)
$`C_{(4)}`$ $`=`$ $`2\stackrel{~}{C}_{(4)}+3C_{(2)}B_{(2)},`$ (207)
for which we define curvatures
$`H_{(3)}`$ $`=`$ $`3B_{(2)},`$ (208)
$`G_{(2n1)}`$ $`=`$ $`(2n1)\{C_{(2n2)}\frac{1}{2}(2n2)(2n3)C_{(2n4)}B_{(2)}\}.`$ (209)
In (209) $`n`$ takes on the values $`n=1,2,3`$, but this will be extended to $`n6`$ below. The corresponding bosonic gauge transformations are
$`\delta B_{(2)}`$ $`=`$ $`\mathrm{\Sigma },`$ (210)
$`\delta C_{(2n2)}`$ $`=`$ $`\mathrm{\Lambda }_{(2n3)}+\frac{1}{2}(2n2)(2n3)\mathrm{\Lambda }_{(2n5)}B_{(2)}.`$ (211)
We now rewrite the supergravity multiplet in these variables:
$`\delta e_\mu ^a`$ $`=`$ $`2i\overline{ฯต}\gamma ^a\psi _\mu `$ (212)
$`\delta \psi _\mu `$ $`=`$ $`D_\mu ฯต\frac{1}{8}\gamma ^{\nu \rho }\sigma _3ฯตH_{\mu \nu \rho }\frac{1}{4}e^\varphi (\gamma G_{(1)})\gamma _\mu (i\sigma _2)ฯต`$ (213)
$`\frac{1}{24}e^\varphi (\gamma G_{(3)})\gamma _\mu \sigma _1ฯต\frac{1}{960}e^\varphi (\gamma G_{(5)})\gamma _\mu (i\sigma _2)ฯต,`$
$`\delta B_{\mu \nu }`$ $`=`$ $`4i\overline{ฯต}\sigma _3\gamma _{[\mu }\psi _{\nu ]},`$ (214)
$`\delta C_{(0)}`$ $`=`$ $`e^\varphi \overline{ฯต}(i\sigma _2)\lambda ,`$ (215)
$`\delta C_{\mu \nu }`$ $`=`$ $`2ie^\varphi \overline{ฯต}\sigma _1\gamma _{[\mu }(\psi _{\nu ]}+\frac{i}{2}\gamma _{\nu ]}\lambda )+C_{(0)}\delta B_{\mu \nu },`$ (216)
$`\delta C_{\mu \nu \rho \sigma }`$ $`=`$ $`4ie^\varphi \overline{ฯต}(i\sigma _2)\gamma _{[\mu \nu \rho }(\psi _{\sigma ]}+\frac{i}{4}\gamma _{\sigma ]}\lambda )+6C_{[\mu \nu }\delta B_{\rho \sigma ]},`$ (217)
$`\delta \lambda `$ $`=`$ $`\frac{i}{2}\gamma ^\mu _\mu \varphi ฯต\frac{i}{24}(\gamma H_{(3)})\sigma _3ฯต`$ (218)
$`+ie^\varphi (\gamma G_{(1)})(i\sigma _2)ฯต+\frac{i}{12}e^\varphi (\gamma G_{(3)})\sigma _1ฯต,`$
$`\delta \varphi `$ $`=`$ $`2\overline{ฯต}\lambda .`$ (219)
The supersymmetry transformations of the RR fields $`C`$ can be summarized as ($`n=1,2,3,๐ซ_n=i\sigma _2`$ for $`n`$ even, $`๐ซ_n=\sigma _1`$ for $`n`$ odd):
$`\delta C_{(2n2)}`$ $`=`$ $`(2n2)ie^\varphi \overline{ฯต}๐ซ_n\gamma _{(2n3)}(\psi _{(1)}+\frac{i}{2n2}\gamma _{(1)}\lambda )`$ (220)
$`+\frac{1}{2}(2n2)(2n3)C_{(2n4)}\delta B_{(2)}.`$
We will now extend this to the higher-rank forms. We define the following RR fields:
$`C_{(6)}`$ $`=`$ $`\frac{1}{4}\stackrel{~}{C}_{(6)}+5C_{(4)}B_{(2)},`$ (221)
$`C_{(8)}`$ $`=`$ $`\frac{1}{2}\stackrel{~}{C}_{(8)}+7C_{(6)}B_{(2)},`$ (222)
$`C_{(10)}`$ $`=`$ $`\frac{3}{4}\stackrel{~}{C}_{(10)}+9C_{(8)}B_{(2)}.`$ (223)
These combinations transform precisely as (220). We have therefore identified the tower of RR forms, in the same form as in . The $`S`$-dual of $`C_{(10)}`$ is however not the field $`B_{(10)}`$ given in . It turns out that $`B_{(10)}`$ corresponds precisely to our $`\stackrel{~}{๐}_{(10)}`$.
The $`S`$-duals of the $`C_{(2n2)}`$ should form a tower of NS-NS forms. If one defines that under $`S`$-duality
$`C_{(2)}iS_{(2)},`$
$`C_{(4)}S_{(4)},`$
$`C_{(6)}iS_{(6)},`$
$`C_{(8)}S_{(8)},`$
$`C_{(10)}iS_{(10)},`$ (224)
then we find
$`S_{(2)}`$ $`=`$ $`\frac{i}{4}\stackrel{~}{B}_{(2)},`$ (225)
$`S_{(4)}`$ $`=`$ $`2\stackrel{~}{C}_{(4)}+6iC_{(2)}S_{(2)},`$ (226)
$`S_{(6)}`$ $`=`$ $`\frac{1}{4}\stackrel{~}{B}_{(6)}+10iC_{(2)}S_{(4)},`$ (227)
$`S_{(8)}`$ $`=`$ $`\frac{1}{2}\stackrel{~}{B}_{(8)}+14iC_{(2)}S_{(6)},`$ (228)
$`S_{(10)}`$ $`=`$ $`\frac{3}{4}\stackrel{~}{B}_{(10)}+18iC_{(2)}S_{(8)}.`$ (229)
For the case $`\mathrm{}=0`$ the supersymmetry variations for $`S_{(n)}`$ are then described by
$`\delta S_{(2n2)}`$ $`=`$ $`(i)^n(2n2)e^{(n2)\varphi }\overline{ฯต}๐ซ_n\gamma _{(2n3)}\left(\psi +\frac{n2}{2n2}\gamma _{(1)}\lambda \right)`$ (231)
$`+i(2n2)(2n3)S_{(2n4)}\delta C_{(2)}`$
where $`S_{(0)}=0`$ and $`๐ซ_n=\sigma ^3`$ for $`n`$ even and $`๐ซ_n=i\sigma ^2`$ for $`n`$ odd.
## 7 Summary and Discussion
In this work we showed that the standard formulation of IIB supergravity can be extended to include a doublet and a quadruplet of ten-form potentials. We argued that no other independent ten-forms can be added to the algebra. We have been using a โ$`SU(1,1)`$-democraticโ formulation, in which all forms are described together with their magnetic duals. Furthermore, all forms transform in a given representation under the duality group $`SL(2,)`$. The previously known RR-ten-form potential $`C_{(10)}`$ is contained in the quadruplet. The other previously known ten-form (named $`B_{(10)}`$ in ) is in the doublet and hence not $`S`$-dual to $`C_{(10)}`$ .
We have shown that all ten-form potentials have a leading term
$$\delta X_{(10)}e^{n\varphi }\overline{ฯต}\gamma _{(9)}\psi \text{at}l=0$$
(232)
in their supersymmetry transformation in string frame where $`X_{(10)}`$ represents a generic ten-form potential.
Such ten-form potentials naturally occur as the leading contribution in Wess-Zumino terms for space-time filling branes with tension $`g_S^n`$. The resulting branes can be found in table 1. These branes and their relevance for theories with sixteen supercharges will be discussed in some detail in a forthcoming paper .
It would be interesting to see how these findings are compatible with the known $`S`$-duality relations between the Heterotic and Type I superstrings. It is well-known that the (Nambu-Goto part of the) tree-level action of the Type I (Heterotic) superstring scales with $`g_S^1`$ ($`g_S^4`$) . The interpretation of the $`g_S^4`$ term at the Heterotic side is not clear. However, the results presented in Table 3, Appendix B, and the $`S`$-duality assignments of the ten-forms (175) open up the possibility to extend this, consistent with $`S`$-duality, to the scaling behaviour $`g_S^1+g_S^3`$ for the Type I superstring and $`g_S^2+g_S^4`$ for the Heterotic superstring such that the Nambu-Goto term at the Heterotic side contains the more conventional $`g_S^2`$ behaviour.
Work on the relation of string- and M-theory with the Kac-Moody algebras $`E_{11}`$ and $`E_{10}`$ has an interesting connection with our results. In it was pointed out that $`E_{10}`$ and $`E_{11}`$ give rise to different IIB ten-form potentials. In particular, $`E_{10}`$ does not give rise to ten-forms, whereas $`E_{11}`$ supports a doublet and a quadruplet of ten-forms . The latter is in agreement with our results.
It will be worthwile to derive the superspace formulation of our results. Note that, although ten-form potentials have identically zero field-strengths, this is not true for the ten-form superpotentials. It would be interesting to calculate the eleven-form curvatures in flat superspace and to see to which kind of Wess-Zumino terms they give rise to. This is the first step towards the construction of a kappa-symmetric Green-Schwarz action for all 9-branes.
## 8 Acknowledgements
We thank Axel Kleinschmidt, Hermann Nicolai and Tomas Ortรญn for useful remarks.
E.B., S.K. and M. de R. are supported by the European Commission FP6 program MRTN-CT-2004-005104 in which E.B., S.K. and M. de R. are associated to Utrecht university. S.K. is supported by a Postdoc-fellowship of the German Academic Exchange Service (DAAD). F.R. is supported by a European Commission Marie Curie Postdoctoral Fellowship, Contract MEIF-CT-2003-500308. The work of E.B. is partially supported by the Spanish grant BFM2003-01090.
## Appendix A Conventions
The Levi-Civita symbol used in this paper is a tensor, and therefore includes the appropriate powers of $`dete`$.
Some useful properties of the complex fermions are:
$`\psi _\mu `$ $`=`$ $`\gamma _{11}\psi _\mu ,`$ (233)
$`\lambda `$ $`=`$ $`\gamma _{11}\lambda ,`$ (234)
$`D_\mu ฯต`$ $`=`$ $`(_\mu +\frac{1}{4}\omega _\mu {}_{}{}^{ab}\gamma _{ab}^{}\frac{i}{2}Q_\mu )ฯต,`$ (235)
$`(\overline{\chi }_1\gamma ^{\mu _1\mathrm{}\mu _n}\chi _2)^{}`$ $`=`$ $`\overline{\chi }_2\gamma ^{\mu _n\mathrm{}\mu _1}\chi _1=(1)^n\overline{\chi }_{1C}\gamma ^{\mu _1\mathrm{}\mu _n}\chi _{2C},`$ (236)
$`\overline{\chi }_1\gamma ^{\mu _1\mathrm{}\mu _n}\chi _2`$ $`=`$ $`(1)^{n(n+1)/2}\overline{\chi }_{2C}\gamma ^{\mu _1\mathrm{}\mu _n}\chi _{1C}.`$ (237)
In these equations $`\chi _i`$ are arbitrary spinors, not necessarily Majorana or Weyl.
For the duality transformations of $`\gamma `$-matrices we have:
$$\gamma ^{\mu _1\mathrm{}\mu _n}=(1)^{{\scriptscriptstyle \frac{1}{2}}n(n1)}\frac{1}{(10n)!}ฯต^{\mu _1\mathrm{}\mu _{10}}\gamma _{\mu _{n+1}\mathrm{}\mu _{10}}\gamma _{11}$$
(238)
The table below gathers the values of the $`U(1)`$ weights of the different fields. The zehnbein $`e_\mu ^a`$ and all form-fields $`A_{(2n)}`$ have weight zero.
## Appendix B Truncations
We briefly sketch how to apply the heterotic and type I truncations to our IIB results and give a list of the fields surviving the truncation.
We first express the complex spinor $`ฯต`$ in terms of two real spinors
$$ฯต=ฯต_1+iฯต_2.$$
(239)
The heterotic truncation is then given by setting
$$ฯต=\pm ฯต_C.$$
(240)
We will work with the โ+โ choice. We also need to make a choice of gauge for the scalars. We make the same choice as in section 6:
$$V_+^2=V_{}^1.$$
(241)
Plugging (240) into the SUSY variation of $`\psi `$ we find
$$\psi =\psi _C.$$
(242)
Similarly, we use the SUSY variations of the other fields to find how the truncation acts on all the fields
$`\psi `$ $`=`$ $`\psi _C,\lambda =\lambda _C,`$ (243)
$`V_+^2`$ $`=`$ $`V_{}^1,V_{}^2=V_+^1,`$ (244)
$`A_{(2)}^1`$ $`=`$ $`A_{(2)}^2,`$ (245)
$`A_{(4)}`$ $`=`$ $`0,`$ (246)
$`A_{(6)}^1`$ $`=`$ $`A_{(6)}^2,`$ (247)
$`A_{(8)}^{11}`$ $`=`$ $`A_{(8)}^{22},A_{(8)}^{12}=0,`$ (248)
$`A_{(10)}^1`$ $`=`$ $`A_{(10)}^2,`$ (249)
$`A_{(10)}^{111}`$ $`=`$ $`A_{(10)}^{222},A_{(10)}^{112}=A_{(10)}^{122}.`$ (250)
We also observe that the relations for the scalars (244) imply, using the reality properties of the scalars and (17), that $`z=\overline{z}`$. This implies that the axion is eliminated by the truncation, using (21) and (25).
The type I truncation is given by setting
$$ฯต=\pm iฯต_C$$
(251)
where we work with the โ+โ-choice again. We choose $`V_+^2=V_{}^1`$ again and find from the SUSY variations
$`ฯต`$ $`=`$ $`iฯต_C,\psi =i\psi _C,\lambda =i\lambda _C,`$ (252)
$`V_+^2`$ $`=`$ $`V_{}^1,V_{}^2=V_+^1`$ (253)
$`A_{(2)}^1`$ $`=`$ $`A_{(2)}^2,`$ (254)
$`A_{(4)}`$ $`=`$ $`0,`$ (255)
$`A_{(6)}^1`$ $`=`$ $`A_{(6)}^2,`$ (256)
$`A_{(8)}^{11}`$ $`=`$ $`A_{(8)}^{22},A_{(8)}^{12}=0,`$ (257)
$`A_{(10)}^1`$ $`=`$ $`A_{(10)}^2,`$ (258)
$`A_{(10)}^{111}`$ $`=`$ $`A_{(10)}^{222},A_{(10)}^{112}=A_{(10)}^{122}.`$ (259)
As in the case of the heterotic truncation, the axion is eliminated by the truncation. We collect the surviving fields of both truncations in table 3.
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# Colourful Simplicial Depth
## 1. Introduction
In statistics there are several measures of the depth of a point $`p`$ in $`^d`$ relative to a fixed set $`S`$ of sample points. Two recent surveys on data depth are \[Alo04\] and \[FR05\], see references therein. The depth measure we are interested in is the simplicial depth of $`p`$, which is the number of simplices generated by points in $`S`$ that contain $`p`$. A point of maximum simplicial depth can be viewed as a type of $`d`$-dimensional median. We would like to obtain a lower bound for the depth of simplicial medians.
To do this, we consider a generalized problem where the sample points are colourful. That is, in dimension $`d`$ we consider sample points given in each of at least $`(d+1)`$ colours. Then we define the colourful simplicial depth of a point $`p`$ relative to this sample to be the number of colourful simplices (i.e. simplices with one vertex of each colour) that contain $`p`$. We focus on the situation where the point $`p`$ is in the intersection of the convex hulls of the individual colours, which is called the core of the configuration.
If $`p`$ is a core point we would typically expect the simplicial depth of $`p`$ to be more than exponential in $`d`$. However, we exhibit configurations where $`p`$ is a core point but is contained in only $`d^2+1`$ colourful simplices. We conjecture that any core point $`p`$ of any $`d`$-dimensional colourful configuration is contained in at least $`d^2+1`$ colourful simplices. Along the way, we notice that both in the colourful and monochrome cases the simplicial depth of points in general position (relative to the sample set) sometimes has pleasant parity properties. We conclude by mentioning some other natural problems relating to the colourful and monochrome simplicial depth.
## 2. Definitions and Background
### 2.1. Simplicial Depth
The (closed) simplicial depth of a point $`p`$ relative to a set $`S`$ of $`n=|S|`$ points in $`^d`$ is the number of (closed) simplices generated by sets of $`(d+1)`$ points from $`S`$ containing $`p`$ in their convex hull. This was introduced by Liu \[Liu90\] as a measure of how representative $`p`$ is of the points in $`S`$. Denote the simplicial depth of $`p`$ relative to $`S`$ as $`\mathrm{depth}_S(p)`$. The simplicial depth of $`p`$ can be interpreted as the probability that $`p`$ lies in a random simplex of $`S`$ times a constant factor of $`n^{d+1}`$ if we sample points from $`S`$ uniformly with replacement, or times $`\left(\genfrac{}{}{0pt}{}{n}{d+1}\right)`$ if we sample without replacement.
We are most interested in the case when $`S\{p\}`$ is in general position, that is for all $`k<d`$ there are no $`k`$-dimensional affine subspaces contain $`k+2`$ points from $`S\{p\}`$. With this assumption, $`p`$ will always be in the interior of any simplices that contain it, so the notions of closed and open simplicial depth coincide. Without this assumption the closed simplicial depth will be larger.
For a set of points $`S`$, define $`f(S)`$ to be the maximum simplicial depth of a point $`p`$ relative to $`S`$, that is:
$$f(S)=\underset{pR^d}{\mathrm{max}}\mathrm{depth}_S(p)$$
(1)
A point $`p`$ maximizing $`f(S)`$ can be understood as a higher dimensional median point. We will call any such $`p`$ a simplicial median. Indeed for $`d=1`$, this is the usual definition of a median $``$. In higher dimensions, this definition retains many desirable properties of the median, such as affine invariance and a high breakdown point (see e.g. \[Alo04\], \[FR05\], \[GSW92\], \[Liu90\]). However, this maximum will not be attained at a point in general position.
We will consider a similar quantity, the maximum simplicial depth of a point $`p`$ that maintains $`S\{p\}`$ in general position:
$$g(S)=\underset{S\{p\}\text{ in general position}}{\mathrm{max}}\mathrm{depth}_S(p)$$
(2)
Equivalently, $`g`$ is the maximum open simplicial depth of a point $`^d`$. In this way the definition of $`g`$ can be extended to the case when $`S`$ is not in general position. While the maximum in (1) will be attained on a discrete set of points in $`^d`$, the maximum in (2) will be attained on an open set. For non-empty $`S`$, we will have $`g(S)<f(S)`$.
### 2.2. Colourful Simplicial Depth
Now consider a situation where points are given in each of $`rd+1`$ colours. Then the sample consists of colourful sets $`S_1,S_2,\mathrm{},S_r`$ which define a colourful configuration $`๐`$. In the following we use a bold font for colourful objects. A colourful simplex from these sets is any simplex whose vertices are chosen from distinct sets. We define $`\mathrm{๐๐๐ฉ๐ญ๐ก}_๐(p)`$ the colourful simplicial depth of $`p`$ relative to the configuration $`๐`$ as the number of colourful simplices containing $`p`$. As with monochrome simplicial depth, colourful depth can be interpreted probabilistically. In the case where $`r=d+1`$, colourful depth corresponds to specifying separate distributions for each vertex of the simplex. Dividing the depth by $`|S_1||S_2|\mathrm{}|S_{d+1}|`$ gives the probability that $`p`$ lies in a random colourful simplex (sampled uniformly).
A choice of sets $`S_1,\mathrm{},S_r`$ specifies a colourful configuration $`๐`$ of points. We call the intersection of the convex hulls of the $`S_i`$โs in a configuration the core of $`๐`$. Bรกrรกny proved that core points are contained in some colourful simplex; this is known at the Colourful Carathรฉodory Theorem \[Bรกr82\]. In the remainder of the paper, except where noted, we assume that all configurations and $`p`$ are in general position and have a non-empty (hence full-dimensional) core. We remark that our results hold under weaker conditions, such as $`p`$ not lying on any hyperplanes generated by points from the configuration.
### 2.3. Background
Even before the notion of simplicial depth was studied in statistics, the question of computing bounds for $`f(S)`$ and $`g(S)`$ given $`n`$ and $`d`$ was studied in the combinatorics and computational geometry communities. The two-dimensional question dates back at least to Kรกrteszi \[Kรกr55\] who showed that for $`n`$ points in the plane, $`g(S)`$ is at most $`\frac{1}{24}(n^3n)`$ for odd $`n`$ and at most $`\frac{1}{24}(n^34n)`$ for even $`n`$, and showed that these bounds were attained when $`S`$ is the set of vertices of a regular $`n`$-gon. In the early 1980โs, Boros and Fรผredi \[BF84\] showed $`g(S)`$ is at least $`\frac{n^3}{27}+O(n^2)`$, and gave configurations achieving this bound.
In a beautiful paper, Bรกrรกny \[Bรกr82\] gave bounds for the monochrome simplicial depth in dimension $`d`$ as an application of his Colourful Carathรฉodory Theorem. He obtained a lower bound by showing that after colouring the points, some point $`p`$ must be contained in many colourful simplices. A key point of Bรกrรกnyโs proof is that a core point $`p`$ of a colourful configuration must lie in at least one colourful simplex. Using this fact, for a set $`S`$ of $`n`$ points in general position in $`^d`$ Bรกrรกny obtains a lower bound of:
$$g(S)\frac{1}{(d+1)^{d+1}}\left(\genfrac{}{}{0pt}{}{n}{d+1}\right)+O(n^d)$$
(3)
This result is asymptotically sharp up to a constant factor as function of $`n`$ (for fixed $`d`$). However, as Bรกrรกny remarks, the constant is probably quite far from the truth. Indeed, he gives a sharp upper bound of:
$$g(S)\frac{1}{2^d(d+1)!}n^{d+1}+O(n^d)$$
(4)
We speculate that the true lower bound is not much less than the upper bound.
One way to improve (3) would be to show that a core point $`p`$ must lie in more than one colorful simplex. In Bรกrรกnyโs original paper, he notes that $`p`$ must in fact lie in at least $`(d+1)`$ colourful simplices, thereby improving (3) to:
$$g(S)\frac{1}{(d+1)^d}\left(\genfrac{}{}{0pt}{}{n}{d+1}\right)+O(n^d)$$
(5)
More generally, if we could show that any core point $`p`$ of a $`d`$-dimensional configuration is contained in at least $`\mu (d)`$ simplices, then we can improve the constant in equation (3) by a factor of $`\mu (d)`$.
## 3. Colourfully Covering the Core
This leads us to ask: What is the minimum number $`\mu (d)`$ of simplices that can contain a core point $`p`$ in a colourful configuration? Given a colourful configuration $`๐`$ with colourful sets $`S_1,\mathrm{},S_r`$ we can define:
$$\mathrm{m}(๐)=\underset{p\mathrm{core}(๐)}{\mathrm{min}}\mathrm{๐๐๐ฉ๐ญ๐ก}_๐(p)$$
(6)
We remark that if $`\mathrm{core}(๐)`$ has a non-empty interior, the minimum in (6) will be attained on an open set of points that are in general position relative to $`๐`$.
In this notation, our objective is to find the minimum value of $`\mathrm{m}(๐)`$ over all configurations $`๐`$ with full-dimensional core in dimension $`d`$. For a fixed $`d`$, it is clear that some configuration with $`(d+1)`$ points in each of $`(d+1)`$ colours attains this minimum, which depends only on the dimension. Hence we can define:
$$\mu (d)=\underset{d\text{ configurations }๐,p\mathrm{core}(๐)}{\mathrm{min}}\mathrm{๐๐๐ฉ๐ญ๐ก}_๐(p)$$
(7)
One might suppose that $`\mathrm{m}(๐)`$ is often large. As a thought experiment, consider choosing a configuration at random. If we take $`(d+1)`$ points in $`^d`$ from a distribution that is โniceโ and centrally symmetric about the origin 0, the probability that $`\mathrm{๐}`$ is contained in their convex hull is $`\frac{1}{2^d}`$ (see e.g. \[WW01\]). This suggests that for random $`๐`$, a typical value for $`\mathrm{๐๐๐ฉ๐ญ๐ก}_๐(\mathrm{๐})`$ would be $`\frac{1}{2^d}(d+1)^{d+1}`$. For a set $`S`$ of $`(d+1)^2`$ points in the plane, plugging this value into Bรกrรกnyโs analysis gives us an estimate of $`g(S)`$ very close to Bรกrรกnyโs upper bound (4). However, it is not immediately clear if we should expect $`\mathrm{m}(๐)`$ to be much smaller than $`\mathrm{๐๐๐ฉ๐ญ๐ก}_๐(\mathrm{๐})`$.
If we take a configuration $`๐^{\mathrm{}}`$ with $`S_1^{\mathrm{}}`$ given by $`(d+1)`$ points in general position and $`S_1^{\mathrm{}}=S_2^{\mathrm{}}=\mathrm{}=S_{d+1}^{\mathrm{}}`$ we get $`\mathrm{m}(๐^{\mathrm{}})=(d+1)!`$. In Section 3.4 we exhibit a configuration $`๐^{}`$ with $`\mathrm{m}(๐^{})d^2+1`$.
In the remainder of the paper, except where noted, we consider configurations with $`(d+1)`$ points in each of $`(d+1)`$ colours.
### 3.1. Preliminaries
In \[BO97b\], Bรกrรกny and Onn consider the problem of colourful linear programming. This is the algorithmic version of the colourful Carathรฉodory problem: Given a core point $`p`$, how can we find a colourful simplex containing $`p`$? They begin with some preprocessing which is also helpful here.
Take a colourful configuration $`๐`$ of $`(d+1)`$ colourful sets in $`^d`$, $`๐=\{S_1,\mathrm{},S_{d+1}\}`$. Take $`p\mathrm{int}(\mathrm{core}(๐))`$. Without loss of generality we assume that the core point $`p=\mathrm{๐}`$. Given any finite set of points $`T^d`$, scaling the points of $`T`$ does not affect whether $`\mathrm{๐}`$ lies in the convex hull of $`T`$ since the coefficients in a convex combination can themselves be rescaled. This allows us to normalize $`๐`$ by rescaling its points to unit vectors.
Let $`\mathrm{conv}(T)`$ be the convex hull of the points in $`T`$ and $`\mathrm{cone}(T)`$ be the set of non-negative linear combinations of points of $`T`$. A cone is simplicial if it can be generated by a set of $`d`$ linearly independent points in $`^d`$. If $`T^d`$ is a set of points, $`\mathrm{๐}T`$, but $`\mathrm{๐}\mathrm{conv}(T)`$, then $`\mathrm{cone}(T)`$ must contain a non-trivial linear subspace of $`^d`$. A closed, convex cone is called pointed if it does not contain such a subspace, so we summarize this as:
###### Lemma 3.1.
Given any finite set of non-zero points $`T^d`$, $`\mathrm{๐}`$ is in $`\mathrm{conv}(T)`$ if and only if $`\mathrm{cone}(T)`$ is not pointed.
When $`T`$ is a finite set of points on the unit $`d`$-sphere $`๐^d^d`$, Lemma 3.1 is equivalent to saying that $`\mathrm{๐}\mathrm{conv}(T)`$ if and only if $`T`$ is not contained in any open hemisphere of $`๐^d`$. One direction is proved by building a hemisphere from a hyperplane through $`\mathrm{๐}`$ whose normal lies in the interior of $`\mathrm{cone}(T)`$ when this cone is pointed. The other direction is proved by observing that an open hemisphere never contains both a point $`p`$ and its antipode $`p`$.
We would like to put Lemma 3.1 in a form that is convenient for counting how many simplices generated from $`T`$ contain 0. To do this, we find it helpful to think about the antipode of one of the points.
###### Lemma 3.2.
If $`T=\{p_1,p_2,\mathrm{},p_{d+1}\}`$ is a set of non-zero affinely independent points in $`^d`$, $`\mathrm{๐}`$ is in $`\mathrm{conv}(T)`$ if and only if the antipode $`p_{d+1}`$ is in $`\mathrm{cone}(p_1,p_2,\mathrm{},p_d)`$.
###### Proof.
Let $`K=\mathrm{cone}(p_1,p_2,\mathrm{},p_d)`$. Since $`K`$ is a cone generated by $`d`$ linearly independent points in $`^d`$, $`K`$ is simplicial and hence pointed. If $`p_{d+1}K`$, then we can write it as a conic combination of the remaining $`p_i`$, that is: $`p_{d+1}=_{i=1}^da_ip_i`$ for some $`a_1,\mathrm{},a_d0`$. Moving the $`p_{d+1}`$ term to the right hand side of the equation and dividing by $`1+_{i=1}^da_i`$ gives $`\mathrm{๐}`$ as a convex combination of the $`p_i`$โs. If $`p_{d+1}`$ is not in $`K`$, then we can strictly separate $`p_{d+1}`$ from $`K`$ with a hyperplane $`H`$ through 0. Then both $`K`$ and $`p_{d+1}`$ lie strictly on the same side of $`H`$, and the cone generated by $`T`$ must be pointed. โ
### 3.2. A Variational Approach
Take a point $`p`$ from a finite set $`S^d`$. Call a simplex generated by points in $`S`$ a $`p`$-simplex if $`p`$ is one of the points used to generate the simplex, and call a $`p`$-simplex zero-containing if it contains $`\mathrm{๐}`$ in its interior. Define $`z_S(p)`$ to be the number of zero-containing $`p`$-simplices for a given $`S`$.
Lemma 3.2 tells us that $`z_S(p)`$ is the number of simplicial cones generated by $`S\{p\}`$ that contain $`p`$. We find it useful to think about what happens to $`z_S(p)`$ if we move $`p`$ while fixing the remaining points of $`S`$. This is particularly illustrative if we confine $`p`$ to the surface of the unit sphere $`๐^d`$ centred at 0.
Let $`U=S\{p\}`$ with $`|U|=u`$. Initially $`z_S(p)`$ will be the number of simplicial cones generated by sets of $`d`$ points from $`U`$ that contain $`p`$. Now consider what happens as $`p`$ (and hence $`p`$) move. The value of $`z_S(p)`$ will stay fixed until $`p`$ crosses the boundary of some simplicial cone from $`U`$. These boundaries are defined by the hyperplanes generated by $`\mathrm{๐}`$ and sets of $`(d1)`$ points from $`U`$. Taking all $`(d1)`$ sets from $`U`$, we can generate all such boundaries. They divide the surface of $`๐^d`$ into open cells that are $`(d1)`$-dimensional open sets. We can define the depth of a cell of $`S`$ to be the number of simplicial cones generated by $`S`$ containing any given point in the interior of the cell.
Consider moving $`p`$ along the surface of $`๐^d`$ to a new point $`p^{}`$. If $`p`$ and $`p^{}`$ are in the same cell, we will have $`z_S(p)=z_S(p^{})`$. Now suppose $`p`$ is in a cell $`C`$ adjacent to the cell $`C^{}`$ containing $`p^{}`$. Then as we move from $`p`$ to $`p^{}`$ we cross a single hyperplane $`H`$ defined by a set $`U^0`$ of $`(d1)`$ points from $`U`$ belonging to $`H`$. Letโs say that $`p`$ is on the left of $`H`$ and $`p^{}`$ is on the right. For the moment we will assume that only $`(d1)`$ points of $`U`$ lie on $`H`$. Let $`U^{}`$ be the set of $`k`$ points from $`U`$ on the left of $`H`$, and $`U^+`$ be the $`uk(d1)`$ points from $`U`$ on the right. Since $`p`$ is in a cell bordered by $`H`$, it lies in the cone defined by the points from $`U^0`$ and $`x`$ for any point $`xU^{}`$. On the other hand, $`p`$ is separated by $`H`$ from the cones formed by $`U^0`$ and $`y`$ for any $`yU^+`$. Hence $`p`$ is contained in exactly $`k`$ simplicial cones from $`S`$ generated by $`U^0`$ and a single other point. Similarly, $`p^{}`$ is contained in exactly $`uk(d1)`$ such cones. Simplicial cones that do not contain $`U^0`$ in their generating set will not have $`H`$ as a facet, so they will contain $`p`$ if and only if they contain $`p^{}`$. Suppose $`p`$ is in $`l`$ such cones. Then $`z_S(p)=l+k`$, while $`z_S(p^{})=l+uk(d1)`$.
We conclude that given the value of $`z_S(p)`$ at some point $`p`$, we can in principle compute $`z_S(p^{})`$ for any other point $`p^{}`$ by tracing a path from $`p`$ to $`p^{}`$, and seeing how each hyperplane generated from points in $`U=S\{p\}`$ divides the points of $`U`$. To do this formally, we need a topological lemma that says we can always draw a path between two points on $`๐^d`$ that crosses only hyperplanes from $`U`$ (as opposed to passing through cones generated by fewer than $`(d1)`$ points). This reduces to the following fact which can be proved using algebraic topology, see for example \[Mun84\]:
###### Lemma 3.3.
The sphere $`๐^d`$, a $`(d1)`$ dimensional manifold, remains path connected after removing finitely many $`(d3)`$-dimensional manifolds.
### 3.3. Parity
The variational approach to computing $`z_S(p)`$ explains the following parity phenomenon:
###### Proposition 3.4.
For any colourful configuration $`๐`$ of $`(d+1)`$ points in each of $`(d+1)`$ colours in odd dimension $`d`$ and any point $`p`$ with $`๐`$ and $`p`$ in general position, the colourful simplicial depth of $`p`$ with respect to $`๐`$ is even.
The authors were surprised by this fact while experimenting with configurations. However, it is easy to explain this via a colourful version of the method described in Section 3.2. Suppose we begin with a configuration $`๐^0`$ with $`(d+1)`$ points in each of $`(d+1)`$ colours clustered near the North Pole of $`๐^d`$. (We remarked in Section 3.1 that it is sufficient to consider configurations on the surface of $`๐^d)`$. Then we can move one point at a time from its initial position in $`๐^0`$ to its final position in $`๐`$ generating a sequence of configurations $`๐^0,๐^1,๐^2,\mathrm{},๐^{(d+1)^2}=๐`$. Clearly $`\mathrm{๐๐๐ฉ๐ญ๐ก}_{๐^0}(\mathrm{๐})=0`$. As we move a given point $`p_i`$ of colour $`j`$ from its initial position in $`๐^0`$ (and $`๐^i`$) to its final position in $`๐`$ (and $`๐^{i+1}`$), we need only to know what happens when the antipode $`p_i`$ crosses colourful hyperplanes defined by a set of $`(d1)`$ points of $`(d1)`$ colours, and not of colour $`j`$. Such a colourful hyperplane $`H`$ will miss only one other colour, $`j^{}`$. There will be $`k`$ points of colour $`j^{}`$ on one side of $`H`$, and $`(d+1k)`$ on the other side. Here we are assuming that the points from $`๐`$ are in general position, but we can argue by continuity that this assumption is not necessary. As $`p_i`$ crosses $`H`$ the number of simplicial cones containing $`p_i`$ generated by points from $`H`$ and a point of colour $`j^{}`$ changes from $`k`$ to $`(d+1k)`$. As long as $`(d+1)`$ is even, the parity doesnโt change.
Examining this proof, we can see that Proposition 3.4 can be generalized:
###### Theorem 3.5.
If $`๐=\{S_1,S_2,\mathrm{},S_r\}`$ is a $`d`$-dimensional colourful configuration of points and for each $`i=1,2,\mathrm{},r`$ we have $`|S_i|`$ even, and $`p`$ is and any point $`p`$ with $`๐`$ and $`p`$ in general position, then the colourful simplicial depth of $`p`$ with respect to $`๐`$ is also even.
For monochrome depth, as we move point $`p`$ around $`๐^d`$ we need to consider all possible hyperplanes formed from $`S\{p\}`$. Using the same reasoning as Theorem 3.5 we get:
###### Theorem 3.6.
If $`S`$ is a set of $`n`$ points in $`^d`$, and $`nd`$ is even, and $`p`$ is a point such that $`S\{p\}`$ is in general position, then the simplicial depth of $`p`$ with respect to $`S`$ is even.
###### Remark 3.7.
The variational approach suggested in Section 3.2 has appeared in various guises in discussions of monochrome simplicial depth. In particular, it underlies many of the algorithms suggested for computing monochrome simplicial depth. Several such algorithms have been proposed recently, see for example the discussion in \[Alo04\]. Many of these focus on the 2-dimensional problem, \[GSW92\] and \[CO01\] use variational ideas in 3- and 4-dimensional algorithms.
For this reason, we believed that Theorem 3.6 existed as folklore for some time. Baker remarks on the two-dimensional version in a recreational mathematics note \[Bak78\], but this fact, which impressed the authors with its simple elegance, seems curiously neglected in the literature. We speculate that one reason for this is that in statistics the focus has been on computing the depth of the sample points themselves, which are not in general position and do not retain nice parity conditions.
### 3.4. Configurations with Small Minimal Colourful Depth
We now describe how to build a colourful configuration $`๐^{}`$ that contains $`\mathrm{๐}`$ in its core, but where only $`d^2+1`$ colourful simplices contain 0. Our strategy is to fix the first $`d`$ colourful sets $`S_1^{},S_2^{},\mathrm{},S_d^{}`$ and then consider possible placements of $`(d+1)`$ points $`p_1,p_2,\mathrm{},p_{d+1}`$ to form $`S_{d+1}^{}`$. We will place the points from $`S_1^{}S_2^{}\mathrm{}S_d^{}`$ on the sphere $`๐^d`$ in such a way that some regions of $`๐^d`$ are sparsely covered by simplices from $`S_1^{}S_2^{}\mathrm{}S_d^{}`$.
We begin by fixing $`ฯต=\frac{1}{100d}`$. We will place the points from $`๐^{}`$ in three locations on $`๐^d`$. The first on the Tropic of Capricorn, which we define to be the set of points on $`๐^d`$ whose $`d`$th coordinate is $`2ฯต`$. The second is on the Tropic of Cancer, whose $`d`$th coordinate is $`ฯต`$. The two tropics are topologically copies of $`๐^{d1}`$, but unlike their namesakes they are not equally spaced from the equator. The final region is the polar region of points in $`๐^d`$ which are within $`ฯต`$ of the North Pole $`p_{\mathrm{north}}=(0,0,\mathrm{},0,1)`$.
Now letโs fix the positions of the points $`\{x_1,x_2,\mathrm{},x_{d+1}\}S_1^{}`$. Take:
$$x_1=(\sqrt{14ฯต^2},0,0,\mathrm{},0,2ฯต)\text{ and }x_2=(\sqrt{1ฯต^2},0,0,\mathrm{},0,ฯต)$$
Note that the line segment between $`x_1`$ and $`x_2`$ passes just below the origin in the sense that it contains a point whose first $`(d1)`$ coordinates are 0, and whose $`d`$th coordinate is negative (and small). We now place the remaining points $`x_3,\mathrm{},x_{d+1}`$ in the polar region in such a way as to ensure that $`\mathrm{๐}\mathrm{int}(\mathrm{conv}(S_1^{}))`$. For $`d=2`$ we can do this by placing $`x_3`$ at the North Pole. For $`d3`$ we can place the points on the section of the Arctic Circle (points with distance $`ฯต`$ to the North Pole) with zero initial coordinate. Topologically the Arctic Circle is a copy of $`๐^{d2}`$; we can take $`x_3,\mathrm{},x_{d+1}`$ to be the vertices of a regular simplex inscribed on this sphere.
The points of colours $`2,3,\mathrm{},d`$ are chosen similarly. The first points from each of the $`d`$ colours are arranged in a regular simplex on Capricorn. The remaining points in the same relative position to the first point, so that each $`S_i^{}`$ is a rotation of $`S_1^{}`$ around the $`d`$th coordinate axis. In particular, for each $`i=1,2,\mathrm{},d`$, the second point of $`S_i^{}`$ will lie on Cancer and the final $`(d1)`$ points will lie in the polar region.
We finish our construction by considering possible placements of the points $`p_1,\mathrm{},p_{d+1}`$ of $`S_{d+1}^{}`$. We want to place the $`p_i`$โs in such a way that their antipodes (the $`p_i`$โs) are contained in few colourful simplicial cones generated from $`๐^{}`$.
Consider the cell $`C_{\mathrm{south}}`$ defined by colours $`1,\mathrm{},d`$ of $`๐^{}`$ on $`๐^d`$ which contains the South Pole $`p_{\mathrm{south}}=(0,0,\mathrm{},0,1)`$. We claim this is exactly the intersection of $`๐^d`$ with the single colourful simplicial cone $`K_{\mathrm{Cap}}`$ defined by the $`d`$ colourful points on Capricorn. This follows since any other colourful cone is generated by a set of $`d`$ coloured points chosen from Capricorn, Cancer and the northern polar region. Fix such a cone and call these sets $`G_{\mathrm{Cap}}`$, $`G_{\mathrm{Can}}`$ and $`G_{\mathrm{Pole}}`$ and let $`K_G=\mathrm{cone}(G_{\mathrm{Cap}}G_{\mathrm{Can}}G_{\mathrm{Pole}})`$. We assume that we have $`|G_{\mathrm{Cap}}|<d`$. We need to show that $`\mathrm{int}(K_{\mathrm{Cap}})\mathrm{int}(K_G)=\mathrm{}`$. To do this, we find a hyperplane separating $`K_{\mathrm{Cap}}`$ and $`K_G`$. If $`G_{\mathrm{Cap}}=\mathrm{}`$ the hyperplane through the Equator will do. Otherwise, take the colours from $`G_{\mathrm{Cap}}`$ and consider any facet $`F`$ of $`K_{\mathrm{Cap}}`$ containing generators of each of these colours. Then $`F`$ separates $`K_{\mathrm{Cap}}`$ from all the polar points and all the Cancer points of colours from $`\{1,2,\mathrm{},n\}G_{\mathrm{Cap}}`$. (To be absolutely proper, in higher dimension we would have to move Capricorn up towards the equator to ensure the separation of the Cancer points, i.e. we would have to reduce the constant $`2ฯต`$ to $`(1+\delta )ฯต`$ for some $`\delta >0`$.) This completes the proof. We conclude that the cell $`C_{\mathrm{south}}`$ is covered only by the colourful cone $`K_{\mathrm{Cap}}`$ and closely approximates the spherical cap bounded by Capricorn.
It is a good strategy to place the antipodes $`p_i`$ in $`C_{\mathrm{south}}`$. If we do this for all of $`S_{d+1}^{}`$, however, the resulting configuration will not have $`\mathrm{๐}\mathrm{conv}(S_{d+1}^{})`$ ($`S_{d+1}^{}`$ would certainly be contained in an open hemisphere). So we must have at least one antipode, say $`p_1`$ above Capricorn. Indeed, if we place the remaining $`p_i`$ below Capricorn, we would need to have $`p_1`$ above the ring of the antipodes of Capricorn. More precisely, this is the set of points on $`๐^d`$ with final coordinate value exactly $`2ฯต`$. In particular, it is above Cancer.
Let $`A=\{a_1,a_2,\mathrm{},a_d\}`$ be the points from $`S_1^{},S_2^{},\mathrm{},S_d^{}`$ on Capricorn. Similarly, let $`B=\{b_1,b_2,\mathrm{},b_d\}`$ be the points on Cancer. Letโs count how many simplicial cones from $`๐^{}`$ must contain $`p_1`$ if we place $`p_1`$ above Cancer. To do this, we start with $`p_1`$ in $`C_{\mathrm{south}}`$ and then move it above Cancer noting which cell boundaries it crosses as suggested in Section 3.2. This structure of the cell boundaries is a topological question, so we find it convenient to remove the $`p_{\mathrm{south}}`$ and equate $`๐^d`$ with $`^{d1}`$.
With the exception of the single colourful cone that contains $`C_{\mathrm{south}}`$, the colourful simplicial cones generated by $`๐^{}`$ correspond to colourful simplices in $`^{d1}`$. The polar points on $`๐^d`$ will be clustered near the origin in $`^{d1}`$. Let $`A^{}=\{a_1^{},\mathrm{},a_d^{}\}`$ and $`B^{}=\{b_1^{},\mathrm{},b_d^{}\}`$ be the projections of $`A`$ and $`B`$ respectively in $`^{d1}`$. Then $`\mathrm{conv}(A^{})`$ and $`\mathrm{conv}(B^{})`$ form nested simplices which contain the projection of the polar region. The boundaries of the colourful simplicial cones on $`๐^d`$ map to facets of simplices in $`^{d1}`$; both are defined by sets of $`(d1)`$ colourful points. Moving $`p_1`$ from below Capricorn to above Cancer corresponds to moving $`p_1^{}`$ from outside $`\mathrm{conv}(A^{})`$ to inside $`\mathrm{conv}(B^{})`$.
Let us now see what simplicial facets $`p_1^{}`$ must cross to do this. If we keep $`p_1^{}`$ far away from the $`a_i^{}`$ and $`b_i^{}`$โs themselves, we can avoid any facets involving the polar points: These facets involve at most $`(d2)`$ generators from $`A^{}`$ and $`B^{}`$, and hence have ends that are at most $`(d3)`$ dimensional manifolds in $`\mathrm{conv}(A^{})\mathrm{int}(\mathrm{conv}(B^{}))`$. The ends can be avoided by Lemma 3.3.
This still leaves $`d2^{d1}`$ colourful facets defined by choosing $`(d1)`$ colourful points from $`A^{}`$ and $`B^{}`$. We can enumerate them by first choosing an index (colour) to omit and then representing the choices of $`a_i^{}`$โs and $`b_i^{}`$โs by a 0-1 vector of length $`(d1)`$. Letting 0 represent the choice of an $`a_i^{}`$, $`\mathrm{conv}(A^{})`$ is bounded by the facets defined by $`d`$ index choices and a vector of 0โs, while $`\mathrm{conv}(B^{})`$ is bounded by the facets defined by $`d`$ index choices and a vector of 1โs. In fact there are $`2^d`$ colourful simplices generated by $`A^{}`$ and $`B^{}`$, and they are enumerated by 0-1 vectors of length $`d`$. Their facets are enumerated by choosing an index to drop from the enumerating sequence. Therefore the sums of the 0-1 vectors enumerating the facets of a given simplex can differ by at most 1.
Now start with $`p_1^{}`$ outside $`\mathrm{conv}(A^{})`$. To bring $`p_1^{}`$ inside $`\mathrm{conv}(B^{})`$, we must start by bringing it into $`\mathrm{conv}(A^{})`$. This involves crossing some boundary face of $`\mathrm{conv}(A^{})`$, say the one defined by $`a_1,\mathrm{},a_{d1}`$. This is enumerated as $`(d,0,0,\mathrm{},0,0)`$. We can proceed through facets $`(d1,0,0,\mathrm{},0,1)`$, $`(d2,0,0,\mathrm{},0,1,1)`$ until finally we cross $`(1,1,1,\mathrm{},1)`$ into a cell of $`\mathrm{conv}(B^{})`$. This involves $`d`$ facet crossings, which is minimal since at each crossing we can only add a single 1 to the 0-1 part of the enumerating vector.
We claim that as $`p_1^{}`$ crosses each facet, it makes a net gain of $`d1`$ containing simplices. At the first facet, $`(d,0,0,\mathrm{},0,0)`$, $`p_1^{}`$ leaves the single exterior simplex defined by the points $`A^{}`$ projected from Capricorn and enters the $`d`$ simplices defined by $`a_1^{},\mathrm{},a_{d1}^{}`$ and the $`d`$ points of colour $`d`$ other than $`a_d^{}`$. At subsequent facet crossings, the same thing happens for the remaining colours: $`p_1^{}`$ leaves the simplex defined by the crossing facet and $`a_i^{}`$. As $`p_1^{}`$ leaves, it enters the simplices defined by this facet and the $`d`$ remaining points of colour $`i`$. Hence the number of simplices containing $`p_1^{}`$ immediately after crossing into $`\mathrm{conv}(B^{})`$ is exactly $`1+d(d1)`$.
We will now return our attention to $`๐^d`$. Denote by $`C_p`$ the cell containing $`p_1`$ whose projection lies inside $`\mathrm{conv}(B^{})`$. From our construction, $`C_p`$ is a cell above Cancer. We want to claim that in fact it contains some point above the set of antipodes of Capricorn, that is, a point whose antipode is in $`C_{\mathrm{south}}`$. This is a complicated geometric calculation. However, we observe that nothing in our topological argument above changes if we change the constant $`2ฯต`$ in our definition of Capricorn to $`cฯต`$ for any $`c0`$. In particular, the cell $`C_p`$ does not degenerate if we move the antipodes of Capricorn towards Cancer by decreasing $`c`$ to 1. Therefore for some $`c>1`$ (this condition maintains $`\mathrm{๐}\mathrm{int}(\mathrm{conv}(S_i^{}))`$ for $`i=1,\mathrm{},d`$), $`C_p`$ includes some point above the antipodes of Capricorn. Any such $`c`$ and point in $`C_p`$ would be sufficient for our construction. We have used $`c=2`$ for concreteness and take it as an article of faith that this is a small enough for our choice of $`ฯต`$.
The construction can now be completed. Take $`p_2`$ to be the midpoint of shortest spherical segment between Capricorn and $`p_1`$ (which lies below Capricorn). Let $`z<2ฯต`$ be the final coordinate of $`p_2`$ and arrange the remaining points so that $`p_2,p_3,\mathrm{},p_{d+1}`$ form a regular simplex on $`๐^d\{x^d|x_d=z\}`$. Then $`\mathrm{๐}`$ is in the convex hull of the $`p_i`$ (and hence $`S_{d+1}^{}`$). Finally we can calculate $`\mathrm{๐๐๐ฉ๐ญ๐ก}_๐^{}(\mathrm{๐})`$ from the location of the $`p_i`$: $`\mathrm{๐}`$ lies in $`1+d(d1)`$ colourful simplices generated with $`p_1`$ and one simplex each including $`p_2,p_3,\mathrm{},p_{d+1}`$. Hence:
$$\mathrm{๐๐๐ฉ๐ญ๐ก}_๐^{}(\mathrm{๐})=1+d(d1)+d=d^2+1$$
###### Remark 3.8.
There are other nice configurations with $`\mathrm{๐๐๐ฉ๐ญ๐ก}_๐^{}(\mathrm{๐})=d^2+1`$. Consider a configuration $`๐^{}`$ similar to $`๐^{}`$ but with the tropics pushed to the north, taking Cancerโs final coordinate to $`3ฯต`$ and Capricornโs to $`ฯต`$. We can then move each of $`p_1,\mathrm{},p_d`$ across Capricorn and the equator through a single boundary facet. Finally place $`p_{d+1}`$ at the South Pole. Using the same analysis as above, we have $`p_1,\mathrm{},p_d`$ points forming $`1+(d1)`$ simplices containing 0, and $`p_{d+1}`$ forming one such simplex for a total of $`d^2+1`$.
Both $`๐^{}`$ and $`๐^{}`$ have symmetry for the first $`d`$ colours, but not the last one. We can also propose a configuration $`๐^{\prime \prime }`$ with symmetry between all the colours. Follow the recipe for $`๐^{}`$ but place one point of each colour on Cancer and Capricorn and place the remaining points in the polar region. This brings a number of technical difficulties, however. The points will not be in general position, since the tropical hyperplanes include $`(d+1)`$ points. It is also a bit less natural to evenly space $`(d+1)`$ points on copies of $`๐^{d1}`$, indeed for $`d=2`$ this construction does not make sense. When there is a nice way to do this for $`d3`$ this (e.g. 4 points on $`๐^2`$) we may end up with some points being antipodes. This would cause $`\mathrm{๐}`$ to be on the faces of some simplices and increase its colourful simplicial depth. Most of these problems can be fixed by perturbing $`๐^{\prime \prime }`$, but even so $`๐^{\prime \prime }`$ is not well-suited to our proof technique. One might also consider configurations that are not confined to the sphere.
### 3.5. Evaluating $`\mu (d)`$
The configuration $`๐^{}`$ of Section 3.4 satisfies $`\mathrm{m}(๐^{})d^2+1`$ where $`\mathrm{m}(๐)`$ is the minimum colourful simplicial depth of core point defined in (6). We would like to prove that $`\mathrm{m}(๐^{})=d^2+1`$ and in fact that for any colourful configuration $`๐`$ we will have $`\mathrm{m}(๐)d^2+1`$, or equivalently $`\mu (d)d^2+1`$. The second half of this proposition clearly implies the first. We suggest it is also more approachable since we can move the core point of minimum depth to $`\mathrm{๐}`$ during preprocessing, whereas a direct attack on $`\mathrm{m}(๐^{})`$ requires understanding the shape of the core of $`๐^{}`$.
Bรกrรกnyโs original Colourful Carathรฉodory theorem is exactly that $`\mu (d)1`$. He further shows that for any $`๐`$ any colourful point from $`๐`$ is part of some generating set for a colourful simplex containing 0. This immediately yields $`\mu (d)d+1`$. In $`๐^{}`$ we see that $`p_2,p_3,\mathrm{},p_{d+1}`$ all generate a unique colourful simplex containing $`\mathrm{๐}`$. Thus the minimum number of colourful simplices containing $`\mathrm{๐}`$ generated by an arbitrary point in a configuration is 1. To get a stronger lower bound than $`\mu (d)d+1`$ we need to understand some global information about configurations.
###### Lemma 3.9.
Fix the sets $`S_1,\mathrm{},S_d`$ from a colourful configuration $`๐`$ with $`\mathrm{๐}`$ in its core, and consider the cells created on $`๐^d`$ by the colourful simplicial cones from these sets. Then every cell has depth at least 1, and if there is a cell of depth 1 it is unique and all other cells have depth at least $`d`$.
###### Proof.
The fact that every cell has depth at least 1 is equivalent to the fact that every colourful point generates some colourful simplex that contains 0, proved in \[Bรกr82\]. Suppose now that there is a cell $`C`$ of depth 1. Any point exiting $`C`$ through a bounding hyperplane will be exiting some colourful simplex. Since the depth of $`C`$ is 1, this will always be the same simplex. Thus the extreme points of $`C`$ must be a colourful set $`A=\{a_1,\mathrm{},a_d\}`$ with $`a_iS_i`$ generating this simplex. We can puncture $`๐^d`$ at $`pC`$ and project $`๐^d\{p\}`$ into $`^{d1}`$. The $`a_i`$โs project to a set $`A^{}=\{a_1^{},\mathrm{},a_d^{}\}`$ that forms a $`(d1)`$-simplex in $`^{d1}`$. The remaining colourful points project to points in $`\mathrm{conv}(A^{})`$.
Take a point $`q`$ inside $`\mathrm{conv}(A^{})`$. We want to show that $`q`$ is contained in at least $`d`$ colourful simplices in addition to $`\mathrm{conv}(A^{})`$ after projection. To do this, it is sufficient to show that if we take any colourful set $`B^{}=\{b_1^{},\mathrm{},b_d^{}\}`$ of projected points with $`b_i^{}`$ of colour $`i`$ and $`A^{}B^{}=\mathrm{}`$, then $`q`$ is in some colourful simplex generated from points of $`A^{}B^{}`$ with some generators from $`B^{}`$. Equivalently, we want to show that $`\mathrm{conv}(A^{})`$ is covered by colourful simplices generated from $`A^{}B^{}`$ (excluding $`\mathrm{conv}(A^{})`$ itself from the covering). Then by partitioning the projections of the colourful points into $`(d+1)`$ colourful sets $`A^{},B_1^{},B_2^{},\mathrm{},B_d^{}`$ we get Lemma 3.9.
Consider the collection X of colourful $`(d1)`$-simplices generated by $`A^{}`$ and $`B^{}`$ in $`^{d1}`$ and let $`\stackrel{~}{X}`$ be the set of points contained in the colourful simplices of $`X`$ other than $`\mathrm{conv}(A^{})`$. The elements of X are all the simplices formed by taking for each colour $`i=1,2,\mathrm{},d`$ either $`a_i^{}`$ or $`b_i^{}`$ as a generating vertex. This construction resembles the $`d`$-dimensional cross-polytope $`\beta _d`$ (the dual of the $`d`$-cube), a regular polytope in $`^d`$ with $`2d`$ vertices and $`2^d`$ facets. The cross-polytope $`\beta _d`$ is generated by taking as vertices the standard unit vectors $`E^+=\{e_1,\mathrm{},e_d\}`$ and their negatives $`E^{}=\{e_1,\mathrm{},e_d\}`$. The facets of $`\beta _d`$ are the convex hulls generated by choosing for each $`i=1,\mathrm{},d`$ either $`e_i`$ or $`e_i`$.
We can see that X is obtained from $`\beta _d`$ as follows: We have $`A^{}B^{}^{d1}`$. Embed $`^{d1}`$ as an affine subspace $`\mathrm{Aff}(A^{})`$ in $`^d`$. Take H to be an affine hyperplane in $`^d`$ parallel to $`\mathrm{Aff}(A^{})`$. For $`i=1,\mathrm{},d`$ let $`p_i`$ be the intersection point of H with the line through $`b_i^{}`$ perpendicular to $`^{d1}`$. Let $`P=\{p_1,\mathrm{}p_d\}`$ and generate a set $`Q`$ of $`(d1)`$-simplices by taking for each $`i=1,\mathrm{},d`$ either $`a_i^{}`$ or $`p_i`$. By construction $`X`$ is the projection of $`Q`$ into $`\mathrm{Aff}(A^{})`$. Now we claim that $`Q`$ is a continuous image of the facets of $`\beta _d`$. We can exhibit such a map by first finding an affine transformation $`T_1`$ with $`T_1(e_i)=a_i^{}`$ for $`i=1,\mathrm{},d`$ and $`T_1(\mathrm{Aff}(E^{}))=\text{H}`$. Note that $`T_1(\beta _d)`$ is a polytope. Then applying a further affine transformation $`t_2`$ to H with $`t_2(e_i)=p_i`$ for $`i=1,\mathrm{},d`$ and extending this to $`T_2`$ on $`^d`$ so that $`T_2`$ fixes $`\mathrm{Aff}(A^{})`$, we see that the composition $`T_2T_1`$ is the required map.
We proceed by contradiction. Assume that $`\stackrel{~}{X}`$ does not cover $`\mathrm{conv}(A^{})`$. Then we can find a retraction of $`\stackrel{~}{X}`$ to its boundary $`(\mathrm{conv}(A^{}))`$. By composing $`T_2`$, the projection taking $`Q`$ onto $`X`$ and the retraction of $`\stackrel{~}{X}`$, we get a retraction of $`T_1(\beta _d)\mathrm{conv}(A^{})`$ onto $`(\mathrm{conv}(A^{}))`$. However, $`T_1((\beta _d))`$ is a $`d`$-dimensional polytope topologically equivalent to $`๐^d`$ and hence $`T_1((\beta _d))\mathrm{conv}(A^{})`$ is topologically equivalent to a $`(d1)`$-dimensional disk $`๐น^{d1}`$. But $`(\mathrm{conv}(A^{}))`$ is topologically equivalent to $`๐^{d1}`$ and a well-known theorem of algebraic topology says that there does not exist a retraction of $`๐น^{d1}`$ to $`๐^{d1}`$ (see for example section 21 of \[Mun84\]). This is the required contradiction, hence the colourful simplices of $`\text{X}\mathrm{conv}(A^{})`$ cover $`\mathrm{conv}(A^{})`$. โ
###### Corollary 3.10.
The minimum colourful simplicial depth of any core point in any colourful configuration is at least $`2d`$. That is, we have $`\mu (d)2d`$.
###### Proof.
It suffices to prove this for a configuration $`๐`$ with $`(d+1)`$ in $`(d+1)`$ colours. Observe that if we have no cell of depth 1 then each of the $`(d+1)`$ points of $`S_{d+1}`$ will generate at least two colourful simplices containing 0, and if we do have such a cell $`C`$, we must place at least one point, say $`p_1S_{d+1}`$ outside of $`C`$ to get $`\mathrm{๐}\mathrm{conv}(S_{d+1})`$. Then $`p_1`$ is generates at least $`d`$ simplices containing $`\mathrm{๐}`$ in addition to the $`d`$ required of the remaining points in $`S_{d+1}`$. โ
### 3.6. The Two-dimensional Case
We will briefly illustrate our methods by describing how core points can be contained in configurations in $`^2`$. Consider such a configuration $`๐=\{X,Y,Z\}`$ with core point $`p`$. We assume general position, and as discussed in Section 3.1, we may without loss of generality take the core point $`p=\mathrm{๐}`$ and place the points of $`๐`$ on the unit circle $`๐^2`$.
Then the points of $`X`$ and $`Y`$ divide $`๐^2`$ into six segments. Let $`X=\{x_1,x_2,x_3\}`$, $`Y=\{y_1,y_2,y_3\}`$. These points generate 9 simplicial cones and divide $`๐^2`$ into 6 segments. The boundaries between cones are simply the rays through the $`x_i`$โs and $`y_i`$โs. Because no three points of $`X`$ or $`Y`$ lie in the same half-circle, each hyperplane through $`\mathrm{๐}`$ and $`x_i`$ divides the $`y_i`$โs 2 to 1 and vice-versa. Then as the antipode of a point from $`Z`$ crosses $`x_i`$ or $`y_i`$ the number of containing simplicial cones changes by exactly one.
To get a configuration $`๐^{}`$ where only 5 simplices contain 0, we take $`x_1=(\sqrt{14ฯต^2},2ฯต)`$, $`x_2=(\sqrt{1ฯต^2},ฯต)`$, $`x_3=(ฯต,\sqrt{1ฯต^2})`$, $`y_1=(\sqrt{14ฯต^2},2ฯต)`$, $`y_2=(\sqrt{1ฯต^2},ฯต)`$, and $`y_3=(ฯต,\sqrt{1ฯต^2})`$. Observe there is a large cell of depth 1 between $`x_1`$ and $`y_1`$. The reader can verify that the sequence of colourful cell depths is: 1,2,3,4,3,2.
Let $`Z=\{z_1,z_2,z_3\}`$. Place $`z_2=(\sqrt{19ฯต^2},3ฯต)`$ and $`z_3=(\sqrt{19ฯต^2},3ฯต)`$ so that their antipodes lie between $`x_1`$ and $`y_1`$. They each generate one simplex containing 0. Finally, to ensure that $`\mathrm{๐}\mathrm{conv}(Z)`$, we see that $`z_1`$ must lie above $`y_2`$ and $`x_2`$. Take $`z_1=(\sqrt{116ฯต^2},4ฯต)`$. Then $`z_1`$ is contained in 3 colourful simplicial cones generated by $`X`$ and $`Y`$. This configuration has $`\mathrm{๐}`$ in the interior of its core and $`\mathrm{๐}`$ lies in $`1+1+3=5`$ colourful simplices.
Using the analysis in Section 3.5 we see that the cells generated by $`X`$ and $`Y`$ have colourful covering depth at least 1. If no cell attains this, then our configuration must yield at least 6 colourful simplices containing 0. If some cell has depth 1, we can place at most two of the $`z_i`$โs in this cell. The remaining $`z_i`$ must then have depth at least 2, for a minimum of 4. In fact, we can strengthen this to show that our configuration is minimal by observing that we cannot place all of $`Z`$ in two adjacent cells. We conclude that $`\mu (2)=5`$. A similar observation in three dimensions shows that $`\mu (3)8`$. Given the construction of Section 3.4 and Proposition 3.4 we know that $`\mu (3)`$ is either 8 or 10. Bรกrรกny and Matouลกek \[BM05\] have shown that $`\mu (3)9`$ which combined with Proposition 3.4 implies that $`\mu (3)=10`$.
## 4. Conclusions
Let us return to our original goals. Using the bound $`\mu (d)2d`$ from Section 3.5, we see that we can improve Bรกrรกnyโs lower bound (5) for the depth of the monochrome simplicial median to:
$$g(S)\frac{2d}{(d+1)^{d+1}}\left(\genfrac{}{}{0pt}{}{n}{d+1}\right)+O(n^d)$$
(8)
This is a modest improvement. Unfortunately, the construction in Section 3.4 shows that simply bounding $`\mu (d)`$ cannot give a stronger bound than:
$$g(S)\frac{d^2+1}{(d+1)^{d+1}}\left(\genfrac{}{}{0pt}{}{n}{d+1}\right)+O(n^d)$$
(9)
Quite recently, Wagner proved exactly the bound (9) in his thesis \[Wag03\] as a special case of his First Selection Lemma. This is, to our knowledge, the first improvement of (5) since Bรกrรกnyโs original paper \[Bรกr82\]. Wagnerโs result uses a continuous version of the Upper Bound Theorem for polytopes and other techniques from probability without any reference to colouring. We find the appearance of the constant $`d^2+1`$, which for us arrives from colourful combinatorics, quite remarkable.
### 4.1. Bounds for Core Point Depth
Recalling that $`\mathrm{m}(๐)`$ is the minimum value of a core point in a configuration $`๐`$ and that $`\mu (d)`$ is the minimum value of $`\mathrm{m}(๐)`$ over all $`d`$-dimensional colourful configurations $`๐`$, our main result is:
###### Theorem 4.1.
The minimal colourful simplicial depth of any interior core point in any colourful configuration is between $`2d`$ and $`d^2+1`$. That is, we have: $`2d\mu (d)d^2+1`$.
###### Conjecture 4.2.
The minimum colourful simplicial depth of any interior core point in any colourful configuration is $`d^2+1`$. That is, we have $`\mu (d)=d^2+1`$.
This conjecture implies that the configuration $`๐^{}`$ minimizes $`\mathrm{m}(๐)`$ for $`d`$-dimensional colourful configurations. It would also give an elementary proof of (9). It is easy to see that this holds for $`d=1`$. As we noted in Section 3.6, Conjecture 4.2 holds for $`d=2`$ and $`d=3`$. The non-uniqueness of configurations attaining $`\mathrm{m}(๐)=d^2+1`$ suggests that any such proof cannot be completely trivial but it may be possible to do this through improved bookkeeping. The authors generated random low-dimensional configurations by computer and did not find any counterexamples to Conjecture 4.2.
###### Remark 4.3.
The lower bound for $`\mu (d)`$ was improved very recently independently by Bรกrรกny and Matouลกek \[BM05\] and Stephen and Thomas \[ST05\] to $`\mathrm{max}(3d,{\displaystyle \frac{1}{5}}d(d+1))`$ for $`d>2`$ and $`{\displaystyle \frac{(d+2)^2}{4}}`$ respectively. We know that $`\mu (1)=2`$, $`\mu (2)=5`$ and $`\mu (3)=10`$. Combining the improved bounds with the parity conditions of Proposition 3.4 we have the following bounds on $`\mu (d)`$ for $`d>3`$:
$$12\mu (4)17,16\mu (5)26,18\mu (6)37,22\mu (7)50,$$
and for $`d>7`$:
$$\frac{(d+2)^2}{4}\mu (d)d^2+1.$$
It is also natural to ask what type of colourful configuration has a core point of maximum colourful simplicial depth. For this question to be interesting, we must fix the number and size of the colourful sets. Hence we restrict our attention to $`d`$-configurations with $`(d+1)`$ points in each of $`(d+1)`$ colours. We also require $`p`$ to lie in the interior of the core since moving to the boundary of a simplex increases the depth. We define:
$$\nu (d)=\underset{d\text{ configurations }๐,p\mathrm{int}(\mathrm{core}(๐))}{\mathrm{max}}\mathrm{๐๐๐ฉ๐ญ๐ก}_๐(p)$$
(10)
Our method is well suited to analyzing $`\nu (d)`$ simply by changing our objective to creating deep cells and placing antipodes in them. We remark that $`\nu (1)=2`$. An analysis similar to that of Section 3.6 shows that $`\nu (2)=9`$. The key observation is after placing two sets of three colourful points on the circle, the sequence of cell depths that we obtain is either 1,2,3,4,3,2 or 3,2,3,2,3,2. In the first case we also need to argue that the cells of depth at least 3 cover less than half the circle and that opposite every point of depth 4 is a point of depth 1.
The minimal core depth configuration $`๐^{}`$ used to prove $`\mu (2)=5`$ is topologically unique, so it is interesting to observe that, up to topology, there are two distinct configurations that contain $`\mathrm{๐}`$ in 9 colourful simplices. The first corresponds to the sequence of cell depths 1,2,3,4,3,2 and contains a point $`z_3`$ that generates a unique 0-containing colourful simplex. The second corresponds to the sequence 3,2,3,2,3,2 and is a combinatorially symmetric configuration where each colourful point is in exactly three 0-containing colourful simplices. The configurations are illustrated in Figure 5.
We can build a configuration $`๐^+`$ with $`\mathrm{๐๐๐ฉ๐ญ๐ก}_{๐^+}(\mathrm{๐})=d^{d+1}+1`$ by following the strategy for $`๐^{}`$ but building a deep cell rather than a shallow one. To do this, we place the polar region points of colour $`i`$ close to the geodesic between $`p_{\mathrm{north}}`$ and the point of colour $`i`$ on Cancer. Then $`p_{\mathrm{north}}`$ is contained in every colourful cone generated by points from Cancer and the polar region (in fact these are all the colourful cones containing $`p_{\mathrm{north}}`$). Hence the cell $`C_{\mathrm{north}}`$ containing $`p_{\mathrm{north}}`$ has depth $`d^d`$. By placing the points of $`S_{d+1}^+`$ so that $`d`$ of their antipodes are in $`C_{\mathrm{north}}`$ and the final antipode is at $`p_{\mathrm{south}}`$, we get $`๐^+`$ with $`\mathrm{๐๐๐ฉ๐ญ๐ก}_{๐^+}(\mathrm{๐})=dd^d+1`$. The two-dimensional $`๐^+`$ appears as the left element of Figure 5.
A more symmetric (but similar) construction places one point of each colour at the vertices of a regular simplex, and the remaining points surround the antipode of the same colour.
It follows that $`\nu (d)d^{d+1}+1`$. We conjecture that this bound is tight. As with Conjecture 4.2 a computer search did not turn up any counterexamples.
###### Conjecture 4.4.
The maximum colourful simplicial depth of any point in the interior of the core of any colourful configuration of $`(d+1)`$ points in each of $`(d+1)`$ colours is $`d^{d+1}+1`$. That is, we have $`\nu (d)=d^{d+1}+1`$.
###### Remark 4.5.
For any $`d`$, there exists a colourful configuration $`๐`$ which contains $`\mathrm{๐}`$ in at least $`32\%`$ of its colourful simplices.
A configuration of $`(d+1)`$ points in each of $`(d+1)`$ colours generates $`(d+1)^{d+1}`$ colourful simplices, so Remark 4.5 follows immediately from the construction of $`๐^+`$. The minimum fraction of colourful simplices containing $`\mathrm{๐}`$ from an $`๐^+`$ configuration is $`82/256`$ attained when $`d=3`$.
## 5. Open Questions
We would like to conclude by mentioning that there are many more natural questions relating to colourful and monochrome simplicial depth. The first is:
###### Question 5.1.
What is a typical value of $`\mathrm{m}(๐)`$ for a random configuration $`๐`$ of $`(d+1)`$ points in each of $`(d+1)`$ colours?
In Section 3, we remarked that such random configurations could be expected to have a simplicial depth on the order of $`\frac{1}{2^d}(d+1)^{d+1}`$ at the origin. We also gave a colourful configuration $`๐^{\mathrm{}}`$ that has $`\mathrm{m}(๐^{\mathrm{}})=(d+1)!`$. However, $`๐^{\mathrm{}}`$ is not in general position. Our construction $`๐^{}`$ from Section 3.4 is in general position and has a low value of $`\mathrm{m}(๐^{})`$. It is not clear if this behaviour is typical, i.e. if most configurations have some point $`p`$ near the edge of the core that drags down $`\mathrm{m}(๐)`$, or if our configuration is statistically unlikely. Indeed we can consider the possibility that all configurations in general position have such a point near the edge of the core.
###### Question 5.2.
What is the maximum value of $`\mathrm{m}(๐)`$ for a colourful configuration $`๐`$ of $`(d+1)`$ points in each of $`(d+1)`$ colours? What if $`๐`$ is not assumed to be in general position?
We observe that in fact our construction of a colourful configuration $`๐^{}`$ with $`\mathrm{m}(๐^{})=d^2+1`$ contains points of high colourful simplicial depth, but away from 0. This leads us to consider the colourful analogues of the functions $`f(S)`$ and $`g(S)`$ of Section 2.1. For a colourful configuration $`๐`$, define:
$$๐(๐)=\underset{p^d}{\mathrm{max}}\mathrm{๐๐๐ฉ๐ญ๐ก}_๐(p)\text{ and }๐ (๐)=\underset{p\text{ in general position}}{\mathrm{max}}\mathrm{๐๐๐ฉ๐ญ๐ก}_๐(p)$$
(11)
We focus on the case where we have $`(d+1)`$ colours. It is clear that given the sizes of the colourful sets $`S_1,\mathrm{},S_{d+1}`$ comprising $`๐`$ that the maximum of $`๐(๐)`$ and $`๐ (๐)`$ is $`|S_1|\mathrm{}|S_{d+1}|`$ and is attained by placing the points of each colour at (or near) the vertices of a simplex. If we restrict $`๐`$ to be a configuration of $`(d+1)`$ points in each of $`(d+1)`$ colours and take the maximum over the interior of the core, we get exactly the question of finding $`\nu (d)`$ (Conjecture 4.4). We are also interested in lower bounds for $`๐(๐)`$ and $`๐ (๐)`$.
###### Question 5.3.
For $`d`$-dimensional configurations consisting of $`n`$ points in each of $`(d+1)`$ colours, find lower bounds for $`๐(๐)`$ and $`๐ (๐)`$.
In a survey paper on the colourful Carathรฉodory theorem, Bรกrรกny and Onn \[BO97a\] mention that the results of \[ABFK92\] can be applied to give a lower bound for $`๐ (๐)`$ when $`n`$ is large of the form:
$$๐ (๐)c_d\left(\genfrac{}{}{0pt}{}{n}{d+1}\right)$$
(12)
Unfortunately, the constant $`c_d`$ is doubly exponential in $`d`$ so the bound is only non-trivial if $`ne^{4d^2}`$. In particular, it sheds no light on the $`n=d+1`$ case.
One can also get a lower bound for $`๐ (๐)`$ directly from the Colourful Tverberg Theorem \[ลฝV92\], which is used to derive the results in \[ABFK92\]:
$$๐ (๐)\frac{1}{4}\left(\frac{n}{d+1}+3\right)$$
(13)
This still does not help for $`n=d+1`$, but for small $`n`$ the bound is stronger than (12) and comes with the additional guarantee that colourful simplices involved are disjoint! This suggests that there is much room for improvement.
### 5.1. Monochrome Questions
The authors would also like to mention that they do not know the answers to some fairly basic questions about monochrome simplicial depth. Recall the maximum closed and open depth functions $`f(S)`$ and $`g(S)`$ for a set of points $`S`$ in $`^d`$ defined in Section 2.1.
###### Question 5.4.
Are the points $`p`$ attaining the maximum $`f(S)`$ in (1) always limit points of the set of maxima attaining $`g(S)`$ in (2)?
We feel that a positive answer to this question would provide a further natural justification for studying $`g(S)`$ in place of $`f(S)`$ when the former is more tractable. Similarly, it would be interesting to get conditions on $`S`$ such that $`f(S)`$ is not much larger than $`g(S)`$.
We are also curious about the expected values of $`f(S)`$ and $`g(S)`$:
###### Question 5.5.
Given $`n`$ points in $`^d`$ distributed independently and symmetrically about 0, what is the expected deepest simplicial depth of the resulting configuration? That is, what is the expected depth of the simplicial median of the points?
Wagner and Welzl \[WW01\] give an expression for the expected depth of 0, but $`\mathrm{๐}`$ will not always be the deepest point. Indeed if $`n=d+1`$ the expected simplicial depth of $`\mathrm{๐}`$ will be $`\frac{1}{2^d}`$ while the simplicial median always has depth 1. For fixed $`d`$ the expected depth of $`\mathrm{๐}`$ is $`\frac{1}{2^d}\left(\genfrac{}{}{0pt}{}{n}{d+1}\right)`$ which has the same asymptotic behaviour as Bรกrรกnyโs sharp upper bound (4) for $`g(S)`$. However, when $`n`$ is not much larger than $`(d+1)`$, the gap between the expected depth of $`\mathrm{๐}`$ and Bรกrรกnyโs upper bound is substantial and it is not clear to us where the expected depth of the simplicial median lies.
Bรกrรกnyโs method of proving (3) combined with a solution to Question 5.1 might lend some insight into Question 5.5, but a direct approach would be better.
## 6. Acknowledgments
We would like to thank Imre Bรกrรกny for discussions which triggered this work and the anonymous referees for suggestions which improved the presentation of this paper.
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# Statistics of Extreme Spacings in Determinantal Random Point Processes
## 1 Introduction
Determinantal (a.k.a. fermion) random point processes were introduced in probability theory by Macchi about thirty years ago (Ma1 , Ma2 , DVJ ). In the last ten years the subject has attracted a considerable attention due to its rich connections to Random Matrix Theory, Combinatorics, Representation Theory, Random Growth Models, Number Theory and several other areas of mathematics. We refer the reader to the recent surveys (S1 , L , HKPV ), and research papers on the subject (BG , DE , GY , Jo1 , Jo2 , LS , Lyt , PV , ST1 , ST2 , SY , S3 , S4 , Yoo ).
In this paper we shall consider determinantal random point processes on the real line with the translation-invariant correlation kernel. In other words, a one particle space $`X`$ is given as $`X=^1,`$ and the space of elementary outcomes $`\mathrm{\Omega }`$ consists of the countable, locally finite particle configurations on the real line
$$\mathrm{\Omega }=\{\xi =(x_i)_{i=\mathrm{}}^+\mathrm{}:\mathrm{\#}(x_i[N,N])<+\mathrm{},N>0\},$$
where $`x_i^1,i=0,\pm 1,\pm 2,\mathrm{},`$ and $`\mathrm{\#}(x_i[N,N])=:\mathrm{\#}([N,N])`$ denotes the number of the particles in the interval $`[N,N].`$ Let us denote the set of the non-negative integers by $`๐ต_+^1=\{0,1,2,\mathrm{}\},`$ and the set of the positive integers by $`๐ฉ=\{1,2,\mathrm{}\}.`$ The $`\sigma `$-algebra $``$ of the measurable subsets of $`\mathrm{\Omega }`$ is generated by the cylinder sets $`C_{I_1,I_2,\mathrm{},I_k}^{n_1,n_2,\mathrm{}n_k}=\{\xi :\mathrm{\#}(I_j)=n_j,j=1,\mathrm{},k\},`$ where $`k`$ is an arbitrary positive integer, $`k๐ฉ,`$ $`I_1,\mathrm{},I_k`$ are arbitrary disjoint subintervals of the real line, and $`n_1,n_2,\mathrm{},n_k๐ต_+^1.`$
A probability measure $`๐ซ`$ on the measurable space $`(\mathrm{\Omega },)`$ defines a random point process $`(\mathrm{\Omega },,๐ซ).`$ A random point process is called determinantal if its $`k`$-point correlation functions have determinantal form
$$\rho _k(x_1,\mathrm{},x_k)=det(K(x_i,x_j))_{i,j=1,\mathrm{},k},k=1,2,\mathrm{},$$
(1)
where $`K(x,y)`$ is usually called the correlation kernel of the determinantal random point process. We remind the reader that $`k`$-point correlation functions are defined in such a way that
$$E\underset{l=1}{\overset{k}{}}\mathrm{\#}(I_l)=_{I_1\times \mathrm{}I_k}\rho _k(x_1,\mathrm{},x_k)๐x_1\mathrm{}๐x_k$$
(2)
for the arbitrary disjoint intervals $`I_1,\mathrm{},I_k.`$ Since the r.h.s. of (1) is non-negative, it follows that the correlation kernel $`K(x,y)`$ has non-negative minors. If, in addition, the integral operator $`K:L^2(^1)L^2(^1),(Kf)(x)=_{\mathrm{}}^+\mathrm{}K(x,y)f(y)๐y,`$ is Hermitian, one can conclude that $`K`$ is non-negative definite, i.e. $`Spec(K)[0,+\mathrm{}).`$ In the Hermitian case one can show that the necessary and sufficient condition on $`K`$ to define a determinantal random point field (1) is
$$0K1,$$
(3)
in other words both $`K`$ and $`1K`$ must be non-negative definite operators (S1 , Ma1 ).
In this paper we consider the translation-invariant kernel
$$K(x,y)=g(yx),\text{where}g(x)=_{\mathrm{}}^+\mathrm{}\mathrm{exp}(2\pi ixt)\varphi (t)๐t,$$
(4)
and $`\varphi (t)`$ is an even real-valued integrable function
$$\varphi (t)=\varphi (t),\varphi L^1(R^1).$$
(5)
It follows from (3) that
$$0\varphi (t)1(a.e.).$$
(6)
In addition, we assume that the following technical conditions are satisfied
$`{\displaystyle _{\mathrm{}}^+\mathrm{}}t^2\varphi (t)๐t<+\mathrm{},`$ (7)
$`|g(x)|{\displaystyle \frac{C}{1+|x|^{\frac{1}{2}+ฯต}}},`$ (8)
$`|g^{}(x)|{\displaystyle \frac{C}{1+|x|^{\frac{1}{2}+ฯต}}},`$ (9)
where $`C`$ is a positive constant, and $`ฯต`$ is an arbitrary small positive constant.
Let $`L`$ be a large positive number. Consider a restriction of a configuration $`\xi `$ to the interval $`[0,L].`$ Let us denote the points of $`\xi [0,L]`$ by $`x_1,x_2,\mathrm{},x_{N(L)},`$ where $`N(L)`$ is the cardinality of $`\xi [0,L].`$ We assume that the points in $`\xi [0,L]`$ are ordered: $`x_1<x_2<\mathrm{}<x_{N(L)}.`$ It is a well known (see e.g. S1 ) that with probability 1 no two particles of a determinantal random point process coincide. We are interested to study the nearest spacings $`\theta _i=x_{i+1}x_i,i=1,\mathrm{},N(L)1,`$ between the neighboring particles. Functional Central Limit Theorem for the empirical distribution function of the nearest spacings of particles in $`[0,L]`$ (in the limit $`L\mathrm{}`$) was proven in S2 for $`K(x,y)=\frac{\mathrm{sin}(\pi x)}{\pi x}`$ (i.e. $`\varphi `$ is the indicator of $`[\frac{1}{2},\frac{1}{2}]`$), and for similar kernels arising in Random Matrix Theory. It was shown in S1 that the result could be extended to a quite general class of translation-invariant correlation kernels.
In this paper we study the smallest nearest spacings. Our main result is the following
###### Theorem 1.1
Let $`(\mathrm{\Omega },,๐ซ)`$ be a determinantal random point process on the real line with the translation-invariant correlation kernel $`K(x,y)=g(yx)`$ satisfying (4)-(9). Then the number of the nearest spacings less than $`s/L^{1/3}`$ in the interval $`[0,L]`$ converges in distribution to the Poisson random variable with the mean $`\alpha s^3,`$ in the limit $`L\mathrm{},`$ where
$$\alpha =\frac{1}{3}g(0)g^{\prime \prime }(0)=\frac{4\pi ^2}{3}_{\mathrm{}}^+\mathrm{}\varphi (t)๐t_{\mathrm{}}^+\mathrm{}t^2\varphi (t)๐t.$$
(10)
Let $`\eta (L)=L^{1/3}\mathrm{min}\{\theta _i,i=1,\mathrm{},N(L)1\}.`$ In other words, $`\eta (L)`$ is the smallest nearest spacing in $`[0,L],`$ rescaled by $`L^{1/3}.`$ Theorem 1 immediately implies
###### Theorem 1.2
Let the conditions in Theorem 1 be satisfied. Then
$$\underset{L\mathrm{}}{lim}\mathrm{Pr}(\eta (L)>s)=\mathrm{exp}(\alpha s^3).$$
(11)
The method of the proof of Theorems 1 and 2 relies on the detailed analysis of $`k`$-point correlation and cluster functions of the $`s`$-modified random point process, introduced in S1 , S2 .
The rest of the paper is organized as follows. Point correlation and cluster functions, and $`s`$-modified random point processes are discussed in Section 2. The proofs of Theorems 1 and 2 are given in Section 3.
We will use the notations $`const,const_k,Const,`$ to denote various positive constants throughout this text. The values of these constants may be different in various parts of the paper. We shall use the notation $`f=O(g)`$ if the ratio $`f/g`$ is bounded from above and below by some positive constants, and the notation $`f=o(g)`$ if the ratio $`f/g`$ goes to zero.
## 2 Correlation and Cluster Functions
We start this section by recalling the definition of a $`k`$-point cluster function (sometimes also known as the Ursell factor). For additional information we refer the reader to R , LP , CL , S2 .
Definition. The $`l`$-point cluster function $`r_l(x_1,\mathrm{},x_l),l=1,2,\mathrm{},`$ of a random point field is defined in terms of the point correlation functions by the formula
$$r_l(x_1,\mathrm{},x_l)=\underset{G}{}(1)^{m1}(m1)!\underset{j=1}{\overset{m}{}}\rho _{|G_j|}(\overline{x}(G_j))$$
(12)
where the sum is over all partitions $`G`$ of $`[l]=\{1,2,\mathrm{},l\}`$ into subsets $`G_1,\mathrm{},G_m,m=1,\mathrm{},l`$, and $`\overline{x}(G_j)=\{x_i:iG_j\},|G_j|=\mathrm{\#}(G_j)`$.
The point correlation functions can be expressed in terms of the point cluster functions as
$$\rho _l(x_1,\mathrm{},x_l)=\underset{G}{}\underset{j=1}{\overset{m}{}}r_{|G_j|}(\overline{x}(G_j)).$$
(13)
The reader can observe that (12) is the Mรถbius inversion formula applied to (13). In particular,
$`\rho _1(x)=r_1(x),`$
$`\rho _2(x_1,x_2)=r_2(x_1,x_2)+r_1(x_1)r_1(x_2),`$
$`\rho _3(x_1,x_2,x_3)=r_3(x_1,x_2,x_3)+r_2(x_1,x_2)r_1(x_3)+r_2(x_1,x_3)r_1(x_2)+`$
$`r_2(x_2,x_3)r_1(x_1)+r_1(x_1)r_1(x_2)r_1(x_3).`$
It follows from (13) and (1) that for determinantal random point fields
$$r_l(x_1,\mathrm{},x_l)=(1)^{l1}\underset{\text{cyclic }\sigma S_l}{}K(x_1,x_{\sigma (1)})K(x_2,x_{\sigma (2)})\mathrm{}K(x_l,x_{\sigma (l)}),$$
(14)
where the sum in (14) is over all cyclic permutations. In other words, for determinantal random processes the difference between the formula (14) for the $`l`$-point cluster function and the formula
$$\rho _l(x_1,\mathrm{},x_l)=\underset{\sigma S_l}{}(1)^\sigma K(x_1,x_{\sigma (1)})K(x_2,x_{\sigma (2)})\mathrm{}K(x_l,x_{\sigma (l)})$$
(15)
for the $`l`$-point correlation function is that in (15) the summation is taken over all permutations in $`S_l,`$ and in (14) the summation is over the cyclic permutations only. One can rewrite (14) as
$$r_l(x_1,\mathrm{},x_l)=(1)^{l1}\frac{1}{l}\underset{\sigma s_l}{}K(x_{\sigma (1)},x_{\sigma (2)})K(x_{\sigma (2)},x_{\sigma (3)})\mathrm{}K(x_{\sigma (l)},x_{\sigma (1)}).$$
It follows from (2) that the integral of the $`k`$-point correlation function over the $`k`$-dimensional cube $`[0,L]^k`$ is equal to the $`k`$-th factorial moment of the counting random variable $`\mathrm{\#}(I),I=[0,L],`$ namely
$$E\mathrm{\#}(I)(\mathrm{\#}(I)1)\mathrm{}(\mathrm{\#}(I)k+1)=_{I^k}\rho _k(x_1,\mathrm{},x_k)๐x_1\mathrm{}๐x_k.$$
(16)
The integral of the $`k`$-point cluster function, in turn, can be expressed as a linear combination of the cumulants of $`\mathrm{\#}([0,L]).`$ Namely, let $`C_k(L)`$ denote the $`k`$-th cumulant of $`\mathrm{\#}([0,L]),`$ and $`V_k(L)=_{[0,L]^k}r_k(x_1,x_2,\mathrm{}x_k)๐x_1๐x_2\mathrm{}๐x_k.`$ Then (see e.g. CL , S2 )
$$\underset{n=1}{\overset{\mathrm{}}{}}\frac{C_n(L)}{n!}z^n=\underset{n=1}{\overset{\mathrm{}}{}}\frac{V_n(L)}{n!}(e^z1)^n.$$
(17)
To apply the machinery of the cluster functions to the problem at hand, we consider a so-called $`s`$-modified random point process, which can be constructed in the following way. We start with a random configuration $`\xi =(x_i)_{i=\mathrm{}}^+\mathrm{}`$ from the original random point field, and keep only those points $`x_i`$ for which there is exactly one neighbor to the right within distance $`s`$, i.e. $`x_{i+1}x_is,x_{i+2}x_i>s.`$ The points $`x_i`$ for which this conditions is not satisfied are thrown away. As a result, we obtain a new random configuration $`\xi (s)\xi ,`$ such that $`\xi (s)=\{x_i:x_{i+1}x_is,x_{i+2}x_i>s\},`$ where $`\mathrm{}<x_2<x_1<x_0<x_1<x_2<\mathrm{}`$ are the points of the original configuration $`\xi .`$ The number of spacings less than $`s`$ of the original random point field in the interval $`[0,L]`$ is related to the number of points of the $`s`$-modified random point field in $`[0,L].`$ As will be shown later, for large $`L`$ and $`sL^{1/3},`$ these two counting random variables coincide with probability very close to 1.
Since the moments and the cumulants of the counting random variable $`\mathrm{\#}([0,L])`$ can be expressed in terms of the integrals of point correlation and cluster functions (16), (17), it is essential to be able to calculate the point correlation and cluster functions of the $`s`$-modified random point process. We shall denote the $`k`$\- point correlation and $`k`$-point cluster functions of the modified random process by $`\rho _k(x_1,\mathrm{},x_k;s)`$ and $`r_k(x_1,\mathrm{},x_k;s)`$, correspondingly. It follows from the inclusion-exclusion principle that, provided $`|x_ix_j|>s,1ijk,`$ one obtains
$`\rho _k(x_1,\mathrm{},x_k;s)={\displaystyle \underset{m=0}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{(1)^m}{m!}}{\displaystyle _{x_1}^{x_1+s}}\mathrm{}{\displaystyle _{x_k}^{x_k+s}}{\displaystyle _{I(x_1,\mathrm{},x_k;s)^m}}`$
$`\rho _{2k+m}(x_1,\mathrm{},x_k,y_1,\mathrm{},y_k,z_1,\mathrm{},z_m)dy_1\mathrm{}dy_kdz_1\mathrm{}dz_m,`$ (18)
where $`I(x_1,\mathrm{},x_k;s)=_{i=1}^m[x_i,x_i+s],`$ and $`I(x_1,\mathrm{},x_k;s)^m=I(x_1,\mathrm{},x_k;s)\times \mathrm{}\times I(x_1,\mathrm{},x_k;s)`$ stands for the $`m`$-th fold Cartesian product of $`I(x_1,\mathrm{},x_k;s)`$ (see e.g. S1 , S2 ).
In the determinantal case (1) the formula for the $`k`$-point cluster function of the $`s`$-modified random process has a somewhat similar structure (S1 , S2 ). Provided $`|x_ix_j|>s,1ijk,`$ one has
$`r_k(x_1,\mathrm{},x_k;s)={\displaystyle \underset{m=0}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{(1)^m}{m!}}{\displaystyle _{x_1}^{x_1+s}}\mathrm{}{\displaystyle _{x_k}^{x_k+s}}{\displaystyle _{I(x_1,\mathrm{},x_k;s)^m}}`$
$`\rho _{2k+m}^{trun}(x_1,\mathrm{},x_k,y_1,\mathrm{},y_k,z_1,\mathrm{},z_m)dy_1\mathrm{}dy_kdz_1\mathrm{}dz_m,`$ (19)
where
$`\rho _{2k+m}^{trun}(x_1,\mathrm{},x_k,y_1,\mathrm{},y_k,z_1,\mathrm{},z_m)={\displaystyle \underset{\sigma S_{2k+m}}{\overset{}{}}}(1)^\sigma K(x_1,\sigma (x_1))\mathrm{}K(x_k,\sigma (x_k))\times `$
$`K(y_1,\sigma (y_1))\mathrm{}K(y_k,\sigma (y_k))K(z_1,\sigma (z_1))\mathrm{}K(z_m,\sigma (z_m)),`$ (20)
where the summation in (20) is over the permutations $`\sigma S_{2k+m}`$ satisfying the property A described below (we note that $`\sigma `$ acts on the set of $`2k+m`$ variables $`x_1,\mathrm{},x_k,y_1,\mathrm{},y_k,z_1,\mathrm{}z_m`$).
Property A
Consider an index $`1ik.`$ Define $`X(i)`$ to be the subset of the set of variables $`\{x_1,\mathrm{},x_k,y_1,\mathrm{},y_k,z_1,\mathrm{}z_m\},`$ that consists of $`x_i,y_i,`$ and those of the variables $`z_1,\mathrm{},z_m,`$ that belong to $`[x_i,x_i+s]`$ (we remind the reader that each $`z_l,1lm,`$ belongs to exactly one interval $`[x_j,x_j+s],1jk).`$ Then for any pair of disjoint integers $`1ijk`$ there exists a positive integer $`r=r(i,j)>0,`$ such that $`\sigma ^r(X(i))X(j)\mathrm{}.`$
We would like to bring to the readerโs attention the fact that the relation between (18) and (19) is, in a sense, quite similar to the relation between (15) and (14).
## 3 Proof of the Main Result
The strategy of the proof is the following. We consider the rescaling $`\stackrel{~}{s}=sL^{\frac{1}{3}},`$ (we shall show that the smallest spacings in the interval $`[0,L]`$ are of order $`L^{1/3}).`$ We shall keep $`s`$ fixed as $`L\mathrm{},`$ so $`\stackrel{~}{s}`$ will be proportional to $`L^{\frac{1}{3}}.`$ We are interested in the asymptotics of the integrals
$$V_k(L)=_{[L,L]^k}r_k(x_1,x_2,\mathrm{}x_k;\stackrel{~}{s})๐x_1๐x_2\mathrm{}๐x_k.$$
(21)
We claim that $`lim_L\mathrm{}V_1(L)=\alpha s^3,`$ where $`\alpha `$ is defined in (10), and $`lim_L\mathrm{}V_k(L)=0,`$ for $`k>1.`$
###### Lemma 1
Let $`V_k(L)`$ be defined as in (28), where $`r_k(x_1,x_2,\mathrm{}x_k;\stackrel{~}{s})`$ is the $`k`$-point cluster function of the $`\stackrel{~}{s}`$-modified random point process introduced above and $`\stackrel{~}{s}=sL^{\frac{1}{3}}.`$ Then
$`\underset{L\mathrm{}}{lim}V_k(L)=\{\begin{array}{cc}\alpha s^3& ifk=1,\\ 0& ifk>1,\end{array}`$ (24)
where $`\alpha `$ has been defined in (10).
The result of Lemma 1, combined with (17), implies that the number of the points of the $`\stackrel{~}{s}`$\- modified random process in the interval $`[0,L]`$ converges in distribution to the Poisson law as $`L\mathrm{}.`$
Once Lemma 1 is proven , we shall show that the number of points in $`[0,L]`$ of the original determinantal process that have at least two neighbors to the right within distance $`s/L^{1/3}`$ is zero with probability very close to 1, provided that $`L`$ is large and $`s`$ stays finite.
Proof of Lemma
We start with $`V_1(L).`$ Consider the one-point correlation function (intensity) of the $`\stackrel{~}{s}`$-modified point process $`\rho _1(x;\stackrel{~}{s}).`$
$$\rho _1(x;\stackrel{~}{s})=\underset{m=0}{\overset{+\mathrm{}}{}}\frac{(1)^m}{m!}_x^{x+\stackrel{~}{s}}\mathrm{}_x^{x+\stackrel{~}{s}}\rho _{m+2}(x,y,z_1,\mathrm{},z_m)๐yd^mz.$$
(25)
We claim that in the determinantal case
$$\rho _1(x;\stackrel{~}{s})=_x^{x+\stackrel{~}{s}}\rho _2(x,y)D(x,y;\stackrel{~}{s})๐y,$$
(26)
where $`D(x,y;\stackrel{~}{s})`$ is the Fredholm determinant of the integral operator $`\stackrel{~}{K}`$ on $`L^2([x,x+\stackrel{~}{s}]),`$
$$D(x,y;\stackrel{~}{s})=det(1\stackrel{~}{K}),\stackrel{~}{K}:L^2([x,x+\stackrel{~}{s}])L^2([x,x+\stackrel{~}{s}]).$$
(27)
The kernel $`\stackrel{~}{K}(u,v)`$ in (27), (26) depends on $`x`$ and $`y,`$ and is given by the formula
$`\stackrel{~}{K}(u,v)`$ $`=`$ $`K(u,v)K(u,x)T_{11}(x,y)K(x,v)K(u,x)T_{12}(x,y)K(y,v)`$ (28)
$``$ $`K(u,y)T_{21}(x,y)K(x,v)K(u,y)T_{22}(x,y)K(y,v),`$
where
$$\left(\begin{array}{cc}T_{11}(x,y)& T_{12}(x,y)\\ T_{21}(x,y)& T_{22}(x,y)\end{array}\right)=\left(\begin{array}{cc}K(x,x)& K(x,y)\\ K(y,x)& K(y,y)\end{array}\right)^1.$$
Indeed, let us introduce the notation $`K[x_1,\mathrm{},x_k]:=det(K(x_i,x_j))_{i,j=1,\mathrm{},k}.`$ Then,
$$K[x,y,z_1,\mathrm{},z_m]=K[x,y]\stackrel{~}{K}[z_1,\mathrm{},z_m]=\rho _2(x,y)\stackrel{~}{K}[z_1,\mathrm{},z_m].$$
In other words, the conditional distribution of a determinantal random point process with the correlation kernel $`K`$, given there are two particles at $`x`$ and $`y`$ is again a determinantal random point process (on $`^1\{x,y\})`$ with the kernel $`\stackrel{~}{K}`$ (see e.g. ST2 ). This allows us to rewrite (25) as
$$\rho _1(x;\stackrel{~}{s})=_x^{x+\stackrel{~}{s}}\rho _2(x,y)\left(\underset{m=0}{\overset{+\mathrm{}}{}}\frac{(1)^m}{m!}_{[x,x+\stackrel{~}{s}]^m}\stackrel{~}{K}(z_1,\mathrm{},z_m)d^mz\right)๐y,$$
(29)
and (26) follows.
The intensity $`\rho _1(x;\stackrel{~}{s})`$ is constant (i.e. it does not depend on $`x`$) in the translation-invariant case. To estimate $`\rho _1(x;\stackrel{~}{s})=\rho _1(0;\stackrel{~}{s}),`$ we note that
$`\rho _{m+2}(x,y,z_1,\mathrm{},z_m)=K[x,y,z_1,\mathrm{},z_m]K(x,x)K(y,y)K(z_1,z_1)\mathrm{}K(z_m,z_m)`$
$`g(o)^{m+2},`$ (30)
since the determinant of a non-negative definite matrix is bounded from above by the product of the diagonal entries (for the generalization of this result see Lemma 2 below). Then
$`\rho _1(x;\stackrel{~}{s})={\displaystyle _x^{x+\stackrel{~}{s}}}\rho _2(x,y)๐y{\displaystyle _x^{x+\stackrel{~}{s}}}{\displaystyle _x^{x+\stackrel{~}{s}}}\rho _3(x,y,z_1)๐y๐z_1+`$
$`{\displaystyle \frac{1}{2}}{\displaystyle _x^{x+\stackrel{~}{s}}}{\displaystyle _x^{x+\stackrel{~}{s}}}{\displaystyle _x^{x+\stackrel{~}{s}}}\rho _4(x,y,z_1,z_2)๐y๐z_1๐z_2+`$
$`{\displaystyle \underset{m=3}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{(1)^m}{m!}}{\displaystyle _x^{x+\stackrel{~}{s}}}\mathrm{}{\displaystyle _x^{x+\stackrel{~}{s}}}\rho _{m+2}(x,y,z_1,\mathrm{},z_m)๐yd^mz=`$
$`{\displaystyle _x^{x+\stackrel{~}{s}}}\rho _2(x,y)๐y{\displaystyle _x^{x+\stackrel{~}{s}}}{\displaystyle _x^{x+\stackrel{~}{s}}}\rho _3(x,y,z_1)๐y๐z_1+`$
$`{\displaystyle \frac{1}{2}}{\displaystyle _x^{x+\stackrel{~}{s}}}{\displaystyle _x^{x+\stackrel{~}{s}}}{\displaystyle _x^{x+\stackrel{~}{s}}}\rho _4(x,y,z_1,z_2)๐y๐z_1๐z_2+O(\stackrel{~}{s}^4).`$
To estimate $`\rho _3(x,y,z_1)=K[x,y,z_1]`$ and $`\rho _4(x,y,z_1,z_2)=K[x,y,z_1,z_2]`$ we recall that the point correlation functions are given by the determinants, and subtract the first column from the other columns both in $`K[x,y,z_1]`$ and $`K[x,y,z_1,z_2].`$ Since $`y[x,x+\stackrel{~}{s}],z_i[x,x+\stackrel{~}{s}],i1,`$ and the first derivative of $`g`$ is uniformly bounded, we observe that $`\rho _3(x,y,z_1)=K[x,y,z_1]=O(\stackrel{~}{s}^2),\rho _4(x,y,z_1,z_2)=K[x,y,z_1,z_2]=O(\stackrel{~}{s}^3),`$ and, therefore
$$\rho _1(x;\stackrel{~}{s})=_x^{x+\stackrel{~}{s}}\rho _2(x,y)๐y+O(\stackrel{~}{s}^4)=_0^{\stackrel{~}{s}}(g^2(0)g^2(t))๐t+O(\stackrel{~}{s}^4)=\alpha \stackrel{~}{s}^3+O(\stackrel{~}{s}^4),$$
(31)
where $`\alpha `$ has been defined in (10). It follows from (31) that $`lim_L\mathrm{}V_1(L)=lim_L\mathrm{}_0^L\rho _1(x;\stackrel{~}{s})๐x=lim_L\mathrm{}\rho _1(0;\stackrel{~}{s})L=\alpha s^3.`$
Next, we show that $`lim_L\mathrm{}V_k(L)=0`$ for $`k>1.`$ We remind the reader that $`V_k(L)`$ has been defined as $`V_k(L)=_{[L,L]^k}r_k(x_1,x_2,\mathrm{}x_k;\stackrel{~}{s})๐x_1๐x_2\mathrm{}๐x_k.`$ We start with the case $`k=2.`$ Recall (see (19)) that for $`|x_1x_2|>\stackrel{~}{s}`$
$`r_2(x_1,x_2;\stackrel{~}{s})={\displaystyle \underset{m=0}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{(1)^m}{m!}}{\displaystyle _{x_1}^{x_1+\stackrel{~}{s}}}{\displaystyle _{x_2}^{x_2+\stackrel{~}{s}}}{\displaystyle _{I(x_1,x_2;\stackrel{~}{s})^m}}\rho _{4+m}^{trun}(x_1,x_2,y_1,y_2,z_1,\mathrm{},z_m)`$
$`dy_1dy_2dz_1\mathrm{}dz_m,`$ (32)
where $`\rho _{4+m}^{trun}(x_1,x_2,y_1,y_2,z_1,\mathrm{},z_m)`$ has been defined in (20).
As described in the Property A (right after the formula (20)), in order to define $`\rho _{4+m}^{trun}(x_1,x_2,y_1,y_2,z_1,\mathrm{},z_m)`$ one introduces a partition $`X(1)X(2)=\{x_1,x_2,y_1,y_2,z_1,\mathrm{},z_m\},`$ where $`X(1)`$ consists of $`x_1,y_1,`$ and those of the variables $`z_1,\mathrm{},z_m`$ that belong to $`[x_1,x_1+s],`$ and $`X(2)`$ consists of $`x_2,y_2,`$ and those of the variables $`z_1,\mathrm{},z_m`$ that belong to $`[x_2,x_2+s].`$ Let $`X(1)\{z_1,\mathrm{},z_m\}=\{z_{i_1},\mathrm{},z_{i_l}\},`$ and $`X(2)\{z_1,\mathrm{},z_m\}=\{z_{j_1},\mathrm{},z_{j_{ml}}\}.`$ Then
$`\rho _{4+m}^{trun}(x_1,x_2,y_1,y_2,z_1,\mathrm{},z_m)=K[x_1,y_1,x_2,y_2,z_1,\mathrm{},z_m]`$
$`K[x_1,y_1,z_{i_1},\mathrm{},z_{i_l}]K[x_2,y_2,z_{j_1},\mathrm{},z_{j_{ml}}].`$ (33)
We claim that
$$|\rho _{4+m}^{trun}(x_1,x_2,y_1,y_2,z_1,\mathrm{},z_m)|(m+4)!(C\stackrel{~}{s})^{2+m}\frac{const}{1+|x_1x_2|^{1+ฯต}},$$
(34)
where $`const`$ is a constant that may depend on $`s,`$ and $`C`$ is the constant introduced after the formulas (8), (9).
The factor $`(C\stackrel{~}{s})^{2+m}`$ in (33) follows from the uniform bound on the derivative of $`g,`$ and the fact that the $`m+2`$ variables $`y_1,y_2,z_1,\mathrm{},z_m`$ are within distance $`\stackrel{~}{s}`$ from either $`x_1`$ or $`x_2.`$ In other words, one can subtract the first column in the matrices in $`K[x_1,y_1,z_{i_1},\mathrm{},z_{i_l}]`$ and $`K[x_2,y_2,z_{j_1},\mathrm{},z_{j_{ml}}]`$ from the other columns, and subtract the first and the third column in the matrix in $`K[x_1,y_1,x_2,y_2,z_1,\mathrm{},z_m]`$ from the corresponding columns. Such linear operations do not change the value of the determinants, and the new matrices will contain the terms $`g(uw)g(x_jw),`$ in all columns, except those corresponding to $`x_1`$ and $`x_2,`$ where $`j=1,2,`$ and $`u[x_j,x_j+s].`$ Such terms can be estimated from above by $`\left(\mathrm{max}_{x[x_jw,x_j+\stackrel{~}{s}w]}|g^{}(x)|\right)\stackrel{~}{s}.`$
It follows from the definition that $`\rho _{4+m}^{trun}(x_1,x_2,y_1,y_2,z_1,\mathrm{},z_m)`$ can be written as a sum over at most $`(m+4)!`$ permutations, each term being a product $`m+4`$ factors. As we just showed, $`m+2`$ out of those $`m+4`$ factors can be estimated in absolute value by $`C\stackrel{~}{s}.`$ Moreover, Property A implies that at least two factors in each term must be given either by $`g(x_1v),`$ or by $`g(uv)g(x_1v),`$ or by $`g(x_2u)`$, or by $`g(vu)g(x_2u),`$ where $`u[x_1,x_1+\stackrel{~}{s}],v[x_2,x_2+\stackrel{~}{s}].`$ The inequalities (8), (9) imply that these two factors each contribute an upper bound $`\frac{C}{1+(|x_1x_2|\stackrel{~}{s})^{\frac{1}{2}+ฯต}}`$ and the desired estimate (34) follows.
We recall that we defined above $`I(x_1,x_2;\stackrel{~}{s})=[x_1,x_1+\stackrel{~}{s}][x_2,x_2+\stackrel{~}{s}].`$ Then
$`|{\displaystyle _{x_1}^{x_1+\stackrel{~}{s}}}{\displaystyle _{x_2}^{x_2+\stackrel{~}{s}}}{\displaystyle _{I(x_1,x_2;\stackrel{~}{s})^m}}\rho _{4+m}^{trun}(x_1,x_2,y_1,y_2,z_1,\mathrm{},z_m)๐y_1๐y_2๐z_1\mathrm{}๐z_m|`$
$`(m+4)!(C\stackrel{~}{s})^{2+m}(2\stackrel{~}{s})^{2+m}{\displaystyle \frac{const}{1+|x_1x_2|^{1+ฯต}}},`$ (35)
and
$`|r_2(x_1,x_2;\stackrel{~}{s})|{\displaystyle \underset{m=0}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{1}{m!}}(m+4)!(C\stackrel{~}{s})^{2+m}(2\stackrel{~}{s})^{2+m}{\displaystyle \frac{const}{1+|x_1x_2|^{1+ฯต}}}`$
$`{\displaystyle \frac{const(\stackrel{~}{s})^4}{1+|x_1x_2|^{1+ฯต}}}.`$ (36)
We remind the reader that (36) has been derived for $`|x_1x_2|>\stackrel{~}{s}.`$
Since $`\stackrel{~}{s}=sL^{1/3},`$ it follows from the above estimate that
$$\underset{L\mathrm{}}{lim}_0^L_0^Lr_2(x_1,x_2;\stackrel{~}{s})\chi _D(x_1,x_2)๐x_1๐x_2=\underset{L\mathrm{}}{lim}O(\stackrel{~}{s}^4L)=0,$$
(37)
where $`D=\{(x_1,x_2):|x_1x_2|>sL^{1/3}\}.`$
To estimate $`r_2(x_1,x_2;\stackrel{~}{s})`$ on $`D^c,`$ we note that
$$r_2(x_1,x_2;\stackrel{~}{s})=\rho _2(x_1,x_2;\stackrel{~}{s})\rho _1(x_1;\stackrel{~}{s})\rho _1(x_2;\stackrel{~}{s}).$$
It follows then from (31) that $`\rho _1(x_1;\stackrel{~}{s})=\rho _1(x_2;\stackrel{~}{s})consts^3L^1.`$ To estimate $`\rho _2(x_1,x_2;\stackrel{~}{s})`$ we can assume without loss of generality that $`x_2x_1x_2+\stackrel{~}{s}`$. Then
$`\rho _2(x_1,x_2;\stackrel{~}{s})={\displaystyle \underset{m=0}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{(1)^m}{m!}}{\displaystyle _{x_1}^{x_1+s}}{\displaystyle _{I(x_1,x_2;\stackrel{~}{s})^m}}\rho _{3+m}(x_1,x_2,y_1,z_1,\mathrm{},z_m)`$
$`dy_1dz_1\mathrm{}dz_m,`$ (38)
where $`I(x_1,x_2;\stackrel{~}{s})=[x_1,x_1+\stackrel{~}{s}][x_2,x_2+\stackrel{~}{s}]=[x_2,x_1+\stackrel{~}{s}].`$ Subtracting the first column in $`\rho _{3+m}(x_1,x_2,y_1,z_1,\mathrm{},z_m)=K[x_1,x_2,y_1,z_1,\mathrm{},z_m]`$ from the other columns and using (9), we see that $`\rho _{3+m}(x_1,x_2,y_1,z_1,\mathrm{},z_m)(m+3)!(const\stackrel{~}{s})^{m+2}.`$ Integrating over $`y_1,z_1,\mathrm{},z_m`$ and summing over $`m`$ we obtain
$$\rho _2(x_1,x_2;\stackrel{~}{s})const\stackrel{~}{s}^3,$$
(39)
which implies
$`\underset{L\mathrm{}}{lim}{\displaystyle _0^L}{\displaystyle _0^L}r_2(x_1,x_2;\stackrel{~}{s})\chi _{D^c}(x_1,x_2)๐x_1๐x_2=`$
$`\underset{L\mathrm{}}{lim}{\displaystyle _0^L}{\displaystyle _0^L}\rho _2(x_1,x_2;\stackrel{~}{s})\chi _{D^c}(x_1,x_2)๐x_1๐x_2`$
$`\underset{L\mathrm{}}{lim}{\displaystyle _0^L}{\displaystyle _0^L}\rho _1(x_1;\stackrel{~}{s})\rho _1(x_1;\stackrel{~}{s})\chi _{D^c}(x_1,x_2)๐x_1๐x_2=`$
$`\underset{L\mathrm{}}{lim}\left(O(\stackrel{~}{s}^4)LO(\stackrel{~}{s}^7)L\right)=0.`$ (40)
Combining (37) and (40) one obtaines $`lim_L\mathrm{}V_2(L)=0.`$
The argument in the case of general $`k>2`$ is quite similar. Again, we first estimate $`r_k(x_1,x_2,\mathrm{},x_k;\stackrel{~}{s})`$ on $`D=\{(x_1,x_2,\mathrm{},x_k):|x_ix_j|>\stackrel{~}{s},1ijk\}.`$ We will use formulas (19) and (20). To estimate $`\rho _{2k+m}^{trun}(x_1,\mathrm{},x_k,y_1,\mathrm{},y_k,z_1,\mathrm{},z_m)`$ we consider the partition $`X(1)\mathrm{}X(k)=\{x_1,\mathrm{},x_k,y_1,\mathrm{},y_k,z_1,\mathrm{},z_m\},`$ where $`X(i)=`$
$`\{x_1,\mathrm{},x_k,y_1,\mathrm{},y_k,z_1,\mathrm{},z_m\}[x_i,x_i+\stackrel{~}{s}].`$ Let $`X(i)\{z_1,\mathrm{},z_m\}=\{z_1^{(i)},\mathrm{},z_{n_i}^{(i)}\}.`$ Then $`X(i)=\{x_i,y_i,z_1^{(i)},\mathrm{},z_{n_i}^{(i)}\}.`$
It follows from Property A and the inclusion-exclusion principle that
$$\rho _{2k+m}^{trun}(x_1,\mathrm{},x_k,y_1,\mathrm{},y_k,z_1,\mathrm{},z_m)=\underset{G}{}(1)^j(j1)!\underset{l=1}{\overset{j}{}}K_l,$$
(41)
where the summation is over all partitions $`G=G_1\mathrm{}G_j`$ of $`[k]=\{1,2,\mathrm{},k\},j=1,\mathrm{},k,`$ and $`K_l=K[x_i,y_i,z_1^{(i)},\mathrm{},z_{n_i}^{(i)}:iG_l];`$ in other words, $`K_l`$ depends on the variables from $`_{iG_l}X(i),`$ and it is given by the determinant of the matrix built from the correlation kernel $`K(x,y).`$ We claim that
$`|\rho _{2k+m}^{trun}(x_1,\mathrm{},x_k,y_1,\mathrm{},y_k,z_1,\mathrm{},z_m)|(m+2k)!(C\stackrel{~}{s})^{k+m}\times `$
$`\left({\displaystyle \frac{const}{1+|x_1x_2|^{\frac{1}{2}+ฯต}}}{\displaystyle \frac{const}{1+|x_2x_3|^{\frac{1}{2}+ฯต}}}\mathrm{}{\displaystyle \frac{const}{1+|x_kx_1|^{\frac{1}{2}+ฯต}}}+\mathrm{}\right),`$ (42)
where the summation in the last factor of the r.h.s. of (42) is over all $`(k1)!`$ cyclic permutations (for example, the first term in the sum corresponds to the cyclic permutation $`123\mathrm{}k1).`$ We claim that the estimate (42) follows from (8), (9), (41) and Property A. As in the case $`k=2`$ discussed above, we use the fact that each of the $`k+m`$ variables $`y_1,\mathrm{},y_k,z_1,\mathrm{},z_m`$ lies within distance $`\stackrel{~}{s}`$ from one of the $`x_i`$โs, $`i=1,\mathrm{},k.`$ In each $`K_l=K[x_i,y_i,z_1^{(i)},\mathrm{},z_{n_i}^{(i)}:iG_l]`$ in (41) we subtract for each $`iG_l`$ the column corresponding to $`x_i`$ from the column corresponding to $`y_i`$ and from the other columns corresponding to the variables from $`X(i).`$ These linear operations do not change the values of determinants, and, therefore, do not change the value of $`\rho _{2k+m}^{trun}(x_1,\mathrm{},x_k,y_1,\mathrm{},y_k,z_1,\mathrm{},z_m).`$ Now, according to the Property A, we observe that $`\rho _{2k+m}^{trun}`$ is a sum of at most $`(m+2k)!`$ terms. Each term is a product of $`m+2k`$ factors. Property A assures that each term in the sum can be put into correspondence with a cyclic permutation $`\sigma `$ on the set of $`k`$ variables $`x_1,x_2,\mathrm{},x_k,`$ in such a way that $`k`$ out of $`m+2k`$ terms in the product are of the form $`g(x_{\sigma (i)}v),`$ or $`g(uv)g(x_{\sigma (i)}v),i=1,\mathrm{},k,`$ where $`v[x_{\sigma (i+1)},x_{\sigma (i+1)}+\stackrel{~}{s}],`$ and $`\sigma (k+1)=\sigma (1).`$ The bounds (8), (9) then imply (42) in the same manner as has been shown in the case $`k=2.`$ Therefore,
$`|`$ $`{\displaystyle _{x_1}^{x_1+\stackrel{~}{s}}}\mathrm{}{\displaystyle _{x_k}^{x_k+\stackrel{~}{s}}}{\displaystyle _{I(x_1,\mathrm{},x_k;\stackrel{~}{s})^m}}\rho _{2k+m}^{trun}(x_1,\mathrm{},x_k,y_1,\mathrm{},y_k,z_1,\mathrm{},z_m)๐y_1\mathrm{}๐y_k`$ (43)
$``$ $`dz_1\mathrm{}dz_m|(m+2k)!(C\stackrel{~}{s})^{k+m}(ks)^{k+m}\times `$
$`\left({\displaystyle \frac{const}{1+|x_1x_2|^{\frac{1}{2}+ฯต}}}{\displaystyle \frac{const}{1+|x_2x_3|^{\frac{1}{2}+ฯต}}}\mathrm{}{\displaystyle \frac{const}{1+|x_kx_1|^{\frac{1}{2}+ฯต}}}+\mathrm{}\right),`$
provided $`\stackrel{}{x}=(x_1,\mathrm{},x_k)D,`$ i.e. $`|x_ix_j|>\stackrel{~}{s}`$ for $`ij.`$ Then on $`D`$ we have an estimate
$`|r_k(x_1,\mathrm{},x_k;\stackrel{~}{s})|{\displaystyle \underset{m=0}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{1}{m!}}(2k+m)!(C\stackrel{~}{s})^{k+m}(k\stackrel{~}{s})^{k+m}\times `$
$`\left({\displaystyle \frac{const}{1+|x_1x_2|^{\frac{1}{2}+ฯต}}}\times {\displaystyle \frac{const}{1+|x_2x_3|^{\frac{1}{2}+ฯต}}}\times \mathrm{}{\displaystyle \frac{const}{1+|x_kx_1|^{\frac{1}{2}+ฯต}}}+\mathrm{}\right)`$
$`Const_k(\stackrel{~}{s})^{2k}\left({\displaystyle \frac{const}{1+|x_1x_2|^{\frac{1}{2}+ฯต}}}\times {\displaystyle \frac{const}{1+|x_2x_3|^{\frac{1}{2}+ฯต}}}\times \mathrm{}{\displaystyle \frac{const}{1+|x_kx_1|^{\frac{1}{2}+ฯต}}}+\mathrm{}\right),`$
and
$$|_0^L\mathrm{}_0^Lr_k(x_1,\mathrm{},x_k;\stackrel{~}{s})\chi _D(x_1,\mathrm{},x_k)๐x_1\mathrm{}๐x_k|Const_k(L^{1+(k1)(\frac{1}{2}ฯต)})s^{2k}L^{\frac{2k}{3}}.$$
The last estimate implies
$$\underset{L\mathrm{}}{lim}_0^L\mathrm{}_0^Lr_k(x_1,\mathrm{},x_k;\stackrel{~}{s})\chi _D(x_1,\mathrm{},x_k)๐x_1\mathrm{}๐x_k=0,$$
(44)
for $`k3.`$ Our next goal is to show that
$$\underset{L\mathrm{}}{lim}_0^L\mathrm{}_0^Lr_k(x_1,\mathrm{},x_k;\stackrel{~}{s})\chi _D^c(x_1,\mathrm{},x_k)๐x_1\mathrm{}๐x_k=0.$$
(45)
To estimate $`r_k(x_1,\mathrm{},x_k;\stackrel{~}{s})`$ on $`D^c,`$ we rewrite the formula (12) that expresses the $`k`$-point cluster function in terms of point correlation functions:
$`r_k(x_1,\mathrm{},x_k;\stackrel{~}{s})=\rho _k(x_1,\mathrm{},x_k;\stackrel{~}{s})\rho _1(x_1;\stackrel{~}{s})\rho _{k1}(x_2,x_3,\mathrm{},x_k;\stackrel{~}{s})`$
$`\rho _1(x_2;\stackrel{~}{s})\rho _{k1}(x_1,x_3,\mathrm{},x_k;\stackrel{~}{s})\mathrm{}\rho _1(x_k;\stackrel{~}{s})\rho _{k1}(x_1,x_2,\mathrm{},x_{k1};\stackrel{~}{s})+`$
$`2\rho _2(x_1,x_2;\stackrel{~}{s})\rho _{k2}(x_3,\mathrm{},x_k;\stackrel{~}{s})+2\rho _2(x_1,x_3;\stackrel{~}{s})\rho _{k2}(x_2,x_4,\mathrm{},x_k;\stackrel{~}{s})+`$
$`\mathrm{}\rho _2(x_{k1},x_k;\stackrel{~}{s})\rho _{k2}(x_1,x_2,\mathrm{},x_{k2};\stackrel{~}{s})\mathrm{}`$ (46)
We claim that the integral of each of the terms in (46) over $`[0,L]^mD^c`$ has a zero limit as $`L\mathrm{}.`$ To prove it, we consider an arbitrary term in (46),
$$\rho _{k_1}(x_1,x_2,\mathrm{},x_{k_1};\stackrel{~}{s})\rho _{k_2}(x_{k_1+1},x_{k_1+2},\mathrm{},x_{k_1+k_2};\stackrel{~}{s})\mathrm{}\rho _{k_l}(x_{k_1+\mathrm{}k_{l1}+1},\mathrm{},x_k;\stackrel{~}{s}),$$
(47)
where $`k_1+k_2\mathrm{}+k_l=k,k_i1,i=1,\mathrm{},k.`$ We shall estimate $`\rho _{k_1}(x_1,x_2,\mathrm{},x_{k_1};\stackrel{~}{s}),`$ the other $`l1`$ factors are estimated in the same way.
First assume that none of the variables $`x_1,x_2,\mathrm{},x_{k_1}`$ are within distance $`\stackrel{~}{s}`$ from each other. Then one can clearly estimate $`\rho _{k_1}(x_1,x_2,\mathrm{},x_{k_1};\stackrel{~}{s})`$ from above as
$$\rho _{k_1}(x_1,x_2,\mathrm{},x_{k_1};\stackrel{~}{s})_{x_1}^{x_1+\stackrel{~}{s}}\mathrm{}_{x_{k_1}}^{x_{k_1}+\stackrel{~}{s}}\rho _{2k_1}(x_1,\mathrm{},x_{k_1},y_1,\mathrm{}y_{k_1})๐y_1\mathrm{}๐y_{k_1}.$$
(48)
Now, since $`\rho _{2k_1}(x_1,\mathrm{},x_{k_1},y_1,\mathrm{}y_{k_1})=K[x_1,\mathrm{},x_{k_1},y_1,\mathrm{},y_{k_1}],`$ and
$`K[x_1,\mathrm{},x_{k_1},y_1,\mathrm{},y_{k_1}]`$ is the determinant of a $`(2k_1)`$-dimensional (non-negative definite) real symmetric matrix, we can estimate the determinant from above by the product of the determinants of the $`2\times 2`$ diagonal blocks
$$K[x_1,\mathrm{},x_{k_1},y_1,\mathrm{},y_{k_1}]=K[x_1,y_1,\mathrm{},x_{k_1},y_{k_1}]\underset{i=1}{\overset{k_1}{}}K[x_i,y_i].$$
(49)
The bound (49) follows from the Fischer inequality which we state below as Lemma 2.
###### Lemma 2
Let $`M=\left(\begin{array}{cc}A& B\\ B^{}& C\end{array}\right)`$ be a block matrix, let $`A`$ and $`C`$ be $`n\times n`$ and, respectively, $`m\times m`$ non-negative definite matrices, and $`B`$ be a $`m\times n`$ matrix. Then
$`det\left(\begin{array}{cc}A& B\\ B^{}& C\end{array}\right)detAdetC.`$ (52)
Proof
To prove Lemma 2, it is enough to reduce it to the obvious case $`M=\left(\begin{array}{cc}Id& B\\ B^{}& Id\end{array}\right)`$ by appropriate rotations and dilations in $`C^n`$ and $`C^m`$ (see e.g. ST2 ).
As was shown above (see calculations around formula (31)
$$_{x_i}^{x_i+\stackrel{~}{s}}K[x_i,y_i]๐y_iconst(\stackrel{~}{s})^3,$$
(53)
which then implies that
$$\rho _{k_1}(x_1,x_2,\mathrm{},x_{k_1})const^{k_1}(\stackrel{~}{s})^{3k_1}.$$
(54)
If none of the variables are within $`\stackrel{~}{s}`$ from each other in all factors in (47), We infer from (54) that
$`\rho _{k_1}(x_1,x_2,\mathrm{},x_{k_1};\stackrel{~}{s})\rho _{k_2}(x_{k_1+1},x_{k_1+2},\mathrm{},x_{k_1+k_2};\stackrel{~}{s})\times \mathrm{}`$
$`\rho _{k_l}(x_{k_1+\mathrm{}k_{l1}+1},\mathrm{},x_k;\stackrel{~}{s})Const(\stackrel{~}{s})^{3k}=O(L^k),`$ (55)
and the integral of the l.h.s. of (55) over $`[0,L]^kD^c`$ goes to zero as $`L\mathrm{},`$ since $`vol([0,L]^kD^c)=O(L^{k1}).`$
If some of the variables in $`\rho _{k_1}(x_1,\mathrm{},x_{k_1})`$ are within the distance $`\stackrel{~}{s}`$ from one another, the analysis is quite similar. Let us assume, for example, that $`x_1x_2\mathrm{}x_{k_1},`$ and that $`x_ix_{i+1}x_i+\stackrel{~}{s},i=1,\mathrm{},p,`$ and that the rest of the variables $`x_{p+1},\mathrm{},x_{k_1}`$ are not within the distance $`\stackrel{~}{s}`$ from each other. Then
$$\rho _{k_1}(x_1,x_2,\mathrm{},x_{k_1};\stackrel{~}{s})_{x_{p+1}}^{x_{p+1}+\stackrel{~}{s}}\mathrm{}_{x_{k_1}}^{x_{k_1}+\stackrel{~}{s}}\rho _{2k_1p}(x_1,\mathrm{},x_{k_1},y_{p+1},\mathrm{}y_{k_1})๐y_{p+1}\mathrm{}๐y_{k_1}.$$
(56)
One can write
$`\rho _{2k_1p}(x_1,\mathrm{},x_{k_1},y_{p+1},\mathrm{}y_{k_1})=K[x_1,\mathrm{},x_{k_1},y_{p+1},\mathrm{},y_{k_1}]K[x_1,x_2,\mathrm{},x_{p+1},y_{p+1}]`$
$`\times {\displaystyle \underset{i=p+2}{\overset{k_1}{}}}K[x_i,y_i].`$ (57)
As before,
$$_{x_i}^{x_i+\stackrel{~}{s}}K[x_i,y_i]๐y_iconst(\stackrel{~}{s})^3,i=p+1,\mathrm{},k_1.$$
(58)
As for the term $`K[x_1,\mathrm{},x_{p+1},y_{p+1}],`$ one can substract the first column from all other columns, and obtain
$$K[x_1,\mathrm{},x_{p+1},y_{p+1}]Const_k(\stackrel{~}{s})^{2p+2},$$
(59)
since $`|g(x)g(y)|=O(\stackrel{~}{s^2})`$ for $`0x,y\stackrel{~}{s}`$ (we used the fact that $`g^{}(0)=0`$). Combining (58) and (59), and integrating over the $`y`$โs we obtain
$$\rho _{k_1}(x_1,x_2,\mathrm{},x_{k_1};\stackrel{~}{s})Const(\stackrel{~}{s})^{3k_1p+2}=O(L^{k_1+\frac{p}{3}\frac{2}{3}}).$$
(60)
Note, however, that
$$Vol\{(x_1,\mathrm{},x_{k_1}):x_ix_{i+1}x_i+\stackrel{~}{s},i=1,\mathrm{},p\}[0,L]^{k_1}=O(L^{k_1p}\stackrel{~}{s}^p),$$
(61)
and the product of the right hand sides of (60) and (61) goes to zero.
If there are several factors in (55) for which there are variables within distance $`\stackrel{~}{s}`$ from each other, the analysis is very similar, and we leave the details to the reader. Combining all the estimate together, one concludes the integral of the l.h.s. of (55) over $`[0,L]^mD^c`$ goes to zero as $`ล\mathrm{}.`$ This finishes the proof of Lemma.
The result of Lemma 1 and formula (17) imply that $`lim_L\mathrm{}_{n=1}^+\mathrm{}\frac{C_n(L)}{n!}z^n=\alpha s^3(e^z1),`$ where $`\{C_n(L)\}_{n=1}^+\mathrm{}`$ is the sequence of the cumulants of the counting random variable $`N(L),`$ where $`N(L)`$ is the number of the points of the $`\stackrel{~}{s}`$-modified random pont field in the interval $`[0,L].`$ It follows from the definition of the $`\stackrel{~}{s}`$-modified random pont field that $`N(L)=N_1(L)+N_2(L),`$ where $`N_1(L)`$ counts the number of particles of the original random point field that have exactly one neighbor within distance $`\stackrel{~}{s}`$ to the right, and $`N_2(L)`$ counts the number of particles of the original random point field that have more than one neighbor within distance $`\stackrel{~}{s}`$ to the right. We claim that the probability that $`N_2(L)0`$ is going to zero as $`L\mathrm{}.`$ Specifically, we establish
###### Lemma 3
$$\underset{L\mathrm{}}{lim}EN_2(L)=0,$$
(62)
where $`E`$ denotes the mathematical expectation.
Since $`N_2(L)`$ is a non-negative, integer-valued random variable, (62) implies that $`\mathrm{Pr}(N_2(L)0)0`$ as $`L\mathrm{}.`$
The proof of Lemma 3 is elementary. We use the estimate
$$EN_2(L)_0^L\left(_x^{x+\stackrel{~}{s}}_x^{x+\stackrel{~}{s}}\rho _3(x,y_1,y_2)๐y_1๐y_2\right)๐x.$$
(63)
As before one can show that $`\rho _3(x,y_1,y_2)=K[x,y_1,y_2]=O(\stackrel{~}{s}^4),`$ and thus $`EN_2(L)(\stackrel{~}{s})^6L=o(1).`$
Theorem 1 is proven. Theorem 2 immediately follows from Theorem 1. Indeed, the event $`\{\eta (L)>s\}`$ is exactly the event that there are no nearest spacings smaller than $`sL^{\frac{1}{3}}`$ between the particles in $`[0,L].`$
Acknowledgment Research was supported in part by the NSF Grant DMS-0405864.
## Index
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# The impact of solar and atmospheric parameter uncertainties on the measurement of ๐โโ and ๐ฟ
## 1 Introduction
The results of atmospheric, solar, accelerator and reactor neutrino experiments show that flavour mixing occurs not only in the hadronic sector, as it has been known for long, but in the leptonic sector as well. The experimental results point to two very distinct mass differences<sup>1</sup><sup>1</sup>1A third mass difference, $`\mathrm{\Delta }m_{LSND}^21`$ eV<sup>2</sup>, suggested by the LSND experiment , has not being confirmed yet and will not be considered in this paper., $`\mathrm{\Delta }m_{sol}^28.2\times 10^5`$ eV<sup>2</sup> and $`|\mathrm{\Delta }m_{atm}^2|2.5\times 10^3`$ eV<sup>2</sup>. Only two out of the four parameters of the three-family leptonic mixing matrix $`U_{PMNS}`$ are known: $`\theta _{12}32^{}`$ and $`\theta _{23}45^{}`$. The other two parameters, $`\theta _{13}`$ and $`\delta `$, are still unknown: for the mixing angle $`\theta _{13}`$ direct searches at reactors and three-family global analysis of the experimental data give the upper bound $`\theta _{13}11.5^{}`$, whereas for the leptonic CP-violating phase $`\delta `$ we have no informations whatsoever. Two additional discrete unknowns are the sign of the atmospheric mass difference and the $`\theta _{23}`$-octant (if $`\theta _{23}45^{}`$).
The full understanding of the leptonic mixing matrix constitutes, together with the discrimination of the Dirac/Majorana character and the measure of its absolute mass scale, the main neutrino-physics goal for the next decade. However, strong correlations between $`\theta _{13}`$ and $`\delta `$ and the presence of parametric degeneracies in the ($`\theta _{13},\delta `$) parameter space, -, make the simultaneous measurement of the two variables extremely difficult. Several setups have been proposed to face these problems and perform this task, the first option being Super-Beamโs (of which T2K is the first approved one). New machines have been also proposed, such as the $`\beta `$-Beam or the Neutrino Factory .
In the literature, however, the simultaneous measurement of $`\theta _{13}`$ and $`\delta `$ has been normally studied considering the solar and atmospheric mixing parameters as external quantities fixed to their best fit values (see for example Ref. and refs. therein; see also for some recent papers). This is clearly an approximation that has been adopted to get a first insight on the problems related to the ($`\theta _{13},\delta `$) measurement. However, the experimental uncertainties on these parameters can in principle affect the measurement of the unknowns, and it seems important to perform an analysis that goes beyond the two-parameters fits presented in the literature.
In this paper we therefore study, in a systematic way, the impact that โsolarโ (i.e. $`\theta _{12}`$ and $`\mathrm{\Delta }m_{12}^2`$) and โatmosphericโ (i.e. $`\theta _{23}`$ and $`\mathrm{\Delta }m_{23}^2`$) parameters uncertainties have on the measurement of $`\theta _{13}`$ and $`\delta `$ at three of the many proposed setups. By doing this we want to catch the characteristic features of the inclusion of external parameters uncertainties in a ($`\theta _{13},\delta `$) measurement. A complete six-dimensional fit<sup>2</sup><sup>2</sup>2To which one could add in principle other variables such as the matter parameter or systematic errors, . requires a really hard computing effort. The authors of Refs. obviate this problem marginalizing over all the external parameters and reducing the fit to a two-dimentional one. Our approach, conversely, consists of a series of three-parameters fits (taking $`\theta _{12},\mathrm{\Delta }m_{12}^2,\theta _{23}`$ and $`\mathrm{\Delta }m_{23}^2`$ in turn as the third fitting variable) to be compared with standard two-parameters fits in $`\theta _{13}`$ and $`\delta `$. In this way, we realized that the atmospheric parameters are the external inputs whose uncertainties are more important in the reconstruction of ($`\theta _{13},\delta `$), and that must be better measured in future experiments. We have also tried to compare our results with other methods that have been proposed to deal with external parameter uncertainties in the measurement of $`\theta _{13}`$ and $`\delta `$ such as the inclusion of a covariance matrix in two-parameters $`\chi ^2`$โs or the so-called CP-coverage .
We consider here, as exemplificative setups, three CERN-based facilities:
* the 4 MWatt SPL Super-Beam and a $`\gamma 100`$ $`\beta `$-Beam both aimed at the Frรฉjus tunnel where a 440 kT fiducial volume UNO-like Water ฤerenkov detector could be located with a $`L=130`$ km baseline.
* the CERN-based 50 GeV Neutrino Factory (see Ref. and refs. therein), with two detectors of different characteristics to take advantage of both the โgoldenโ and โsilverโ channels $`\nu _e\nu _\mu ,\nu _\tau `$. The two detectors considered are a 40 kT magnetized iron detector located at $`L=3000`$ km and a 4 kT emulsion cloud chamber located at $`L=732`$ km in the Gran Sasso tunnel.
By comparing the results at these three, very different, facilities, we deduce that the impact of the atmospheric parameters uncertainties is a common problem that future experiments looking for $`\theta _{13}`$ and $`\delta `$ will have to face. Of course, this analysis can be done for any of the different setups proposed in the literature and not analyzed here. Our intention is mainly to address, in this paper, the problem of how uncertainties in the atmospheric and solar parameters affect the measurement of ($`\theta _{13},\delta `$) at setups that have been thoroughly discussed than to present a comprehensive comparison between two- and three-parameters fits at all of the facilities proposed in the literature.
The paper is organized as follows: in Sect. 2 we shortly introduce the three facilities and the neutrino-nucleon cross-section; in Sect. 3 we remind the central values and the uncertainties of solar and atmospheric parameters; in Sect. 5 we review the parametric degeneracies in the measurement of $`\theta _{13}`$ and $`\delta `$ in appearance and disappearance channels; in Sect. 4 we introduce the statistical approach used in the paper; in Sect. 6 we present our results for the measurement of $`\theta _{13}`$ and $`\delta `$ taking into account the uncertainties on solar and atmospheric parameters; in Sect. 7 we show the CP-violation discovery potential of the considered facilities taking into account the uncertainties on atmospheric parameters; in Sect. 8 we eventually draw our conclusions. In App. A we compare our statistical approach with other methods; in App. B we present three-parameters fits for the three considered setups for different choices of the input pair ($`\overline{\theta }_{13},\overline{\delta }`$).
## 2 The experimental setup
In this section we describe, briefly, the three facilities that we will use in the following and we remind the neutrino-nucleon cross-section used throughout the paper.
### 2.1 The $`\beta `$-Beam
The $`\beta `$-Beam concept was first introduced in Ref. . It involves producing a beam of $`\beta `$-unstable heavy ions, accelerating them to some reference energy, and allowing them to decay in the straight section of a storage ring, resulting in a very intense neutrino beam. The chosen ions are <sup>6</sup>He, to produce a pure $`\overline{\nu }_e`$ beam, and <sup>18</sup>Ne, to produce a $`\nu _e`$ beam. We follow the setup proposed in Ref. : the $`\gamma `$ ratio for the two ions has been fixed to $`\gamma (^6\mathrm{He})/\gamma (^{18}\mathrm{Ne})=3/5`$, in order to have both ions circulating in the storage ring at the same time; the $`\gamma `$ value has been fixed to $`\gamma _{{}_{}{}^{18}\mathrm{Ne}}=100`$ (i.e., $`\gamma _{{}_{}{}^{6}\mathrm{He}}=60`$) to tune the neutrino/antineutrino mean energy at the maximum of the $`\nu _e\nu _\mu `$ oscillation probability for the CERN to Frรฉjus baseline. A flux of $`2.9\times 10^{18}`$ $`{}_{}{}^{6}\mathrm{He}`$ decays/year and $`1.1\times 10^{18}`$ $`{}_{}{}^{18}\mathrm{Ne}`$ decays/year is assumed. Fig. 1(left) shows the $`\beta `$-Beam neutrino fluxes computed at $`L=130`$ km, keeping $`m_e0`$ and taking into account the three different decay modes of <sup>18</sup>Ne . The mean energy of the $`\overline{\nu }_e`$, $`\nu _e`$ beams for this setup is 0.23 GeV and 0.37 GeV, respectively. Clearly, energy resolution is very poor at such low energy, given the influence of Fermi motion and other nuclear effects. Therefore, in the following all the sensitivities are computed for a counting experiment with no energy cuts . The $`\beta `$-Beam is a clean environment to produce electron-type neutrinos: the main sources of systematic error are the overall flux normalization (that can be controlled with a near detector), the definition of the fiducial volume of the detector and the neutrino-nucleon cross-sections. Alternative $`\beta `$-Beam proposals can be found in Refs. .
### 2.2 The Super-Beam
A Super-Beam is a conventional neutrino beam with a proton intensity higher than that of existing (or under construction) beams such as K2K , NuMI and the CNGS . With respect to the $`\beta `$-Beam and the Neutrino Factory, neutrino beams of a new design, it has the advantage of a well known technology. On the other hand, the flux composition (with $`\nu _\mu `$ as the main component for a $`\pi ^+`$ focusing, plus a small but unavoidable admixture of $`\overline{\nu }_\mu `$, $`\nu _e`$ and $`\overline{\nu }_e`$) limits its sensitivity to $`\nu _\mu \nu _e`$ oscillations.
We follow the setup proposed in Ref. as a reference: a 2.2 GeV proton beam of 4 MWatt power (the SPL), with neutrino fluxes computed in a full simulation of the beamline in Ref. , assuming a decay tunnel length of 60 m. The corresponding fluxes are shown in Fig. 1(right). Notice that this beam was designed, originally, as the first stage of a would-be Neutrino Factory, and it has not been optimized as a facility to look for $`\nu _\mu \nu _e`$ on its own. Such an optimization has been presented in Ref. . Also in this case, as it was for the $`\beta `$-beam, the main source of systematic error are the poorly known neutrino-nucleon cross-sections, the definition of the fiducial volume in the far detector and the overall normalization of the flux (with the additional problem of new background coming from neutrino species not present in the $`\beta `$-Beam flux). For this setup, also, we consider two tentative values of systematic error: an โoptimisticโ 2% and a โpessimisticโ 5%.
### 2.3 The Neutrino Factory
The Neutrino Factory that we consider consists of a SPL-like Super-Beam and a 50 GeV muon storage ring , with $`2\times 10^{20}`$ muons decaying in the straight section of the storage ring per year. Five years of data taking for each muon polarity is envisaged. Two detectors of different technology are considered: a 40 kT Magnetized Iron Detector (MID) at $`L=2810`$ km; and a 4 kT Emulsion Cloud Chamber (ECC) at $`L=732`$ or $`2810`$ km. This proposal corresponds to the design of a possible CERN-based Neutrino Factory Complex, with detectors located at the Gran Sasso Laboratory (the ECC) and at a second site to be defined (the MID and possibly the ECC). Each one of these detectors is especially optimized to look for a particular signal: the โgoldenโ channel $`\nu _e\nu _\mu `$ for the 40 kT MID, and the โsilverโ channel $`\nu _e\nu _\tau `$ for the 4 kT ECC. The corresponding neutrino fluxes are shown in Fig. 2(left).
The detectors background and systematics for this specific facility have been studied in details in Ref. (the Magnetized Iron Detector) and in Ref. (the Emulsion Cloud Chamber).
### 2.4 The neutrino cross-section
An important source of systematic error is our present poor knowledge of the $`\nu N`$ and $`\overline{\nu }N`$ cross-sections for energies below 1 GeV : either there are very few data (the case of neutrinos) or there are no data at all (the case of antineutrinos). On top of that, the few available data have generally not been taken on the target used in the experiments (either water, iron or lead), and the extrapolation from different nuclei is complicated by nuclear effects that at the considered energies play an important role. For definiteness we show in Fig. 2(right) the cross-sections on water used for the Water ฤerenkov detector throughout the paper . Notice that we also used cross-sections on iron and lead for the MID and the ECC, respectively.
## 3 The leptonic mixing parameters
In Tab. 1 we remind the values of solar and atmospheric sector parameters used in the paper, their present uncertainties and the errors expected after a round of new experiments.
In particular, in the second column of Tab. 1 we report the input values for $`\theta _{12},\theta _{23},\mathrm{\Delta }m_{12}^2`$ and $`\mathrm{\Delta }m_{23}^2`$ used in the paper. They correspond to the present best fit values for solar and atmospheric parameters with the only exception of $`\theta _{23}`$, for which we do not use the present best fit value, $`\theta _{23}=45^{}`$, but $`\theta _{23}=40^{}`$ to make manifest the impact of possible octant degeneracies on the results . Notice that throughout the paper the experimentally measured atmospheric mass difference (whose present best fit value will be labelled as $`\mathrm{\Delta }m_{atm}^2`$) will be fitted with the three-family parameter $`|\mathrm{\Delta }m_{23}^2|=|m_3^2m_2^2|`$ (see for a different convention). For the solar mass difference, on the other hand, we can unambiguously identify the three-family parameter $`\mathrm{\Delta }m_{12}^2=m_2^2m_1^2`$ with the experimentally measured quantity, $`\mathrm{\Delta }m_{sol}^2`$. In the third column of Tab. 1 we report the present uncertainties on each of the parameters. Finally, in the fourth column, we present the uncertainties on solar and atmospheric parameters that are expected to be achieved with ongoing or planned experiments. For an estimate of the reduction of solar parameter uncertainties we refer to Ref. . For an estimate of the reduction of atmospheric parameter uncertainties we refer to the Letter of Intent of the T2K-phase I experiment, . The expected error on $`|\mathrm{\Delta }m_{23}^2|`$ for the central value $`|\mathrm{\Delta }m_{23}^2|=\mathrm{\Delta }m_{atm}^2=2.5\times 10^3`$ eV<sup>2</sup> is a function of the sign of the atmospheric mass difference, something that will not be measured at T2K-I. For this reason we present both spreads specifying the chosen hierarchy, using the results of an analysis yet to appear, .
The T2K-I improved bounds are used to analyse the impact of the expected atmospheric uncertainties in the measurement of ($`\theta _{13},\delta `$) at the $`\beta `$-Beam and the Neutrino Factory<sup>3</sup><sup>3</sup>3An updated detailed computation of the expected errors on atmospheric parameters that can be obtained at this facility is lacking, .. On the other hand, in Sect. 5.2 it will be shown that the $`\nu _\mu `$ disappearance channel at the SPL Super-Beam can improve significantly the present uncertainties on the atmospheric parameters. This measure will therefore be combined with the appearance channel when analysing the impact of the atmospheric uncertainties in the measure of ($`\theta _{13},\delta `$) at the Super-Beam.
## 4 Statistical approach
In this section we describe the statistical approach used in the paper to estimate the impact of uncertainties in the atmospheric and solar parameters in the measurement of ($`\theta _{13},\delta `$).
The obvious approach would be to fit the data in $`(N_\alpha +2)`$ parameters, with $`N_\alpha `$ the number of external parameters that are allowed to vary in a given range (e.g., the solar and atmospheric parameters plus the matter density). This procedure, however, is increasingly time-consuming as the number of parameters to be fitted goes up. For this reason, in order to understand how any single parameter affect the measurement, we perform three-parameters fits in $`\theta _{13},\delta `$ and one of the following parameters in turn: $`\theta _{12}`$, $`\mathrm{\Delta }m_{12}^2`$, $`\theta _{23}`$ and $`\mathrm{\Delta }m_{23}^2`$, each of them allowed to vary uniformly in the ranges of Tab. 1 (no gaussian priors are introduced). The matter density has been considered as a fixed quantity throughout the paper .
To perform the three-parameters fits we have constructed grids for the expected number of charged-current events for each facility, each grid in the two unknowns, $`\theta _{13}`$ and $`\delta `$, plus the two measured parameters in the solar sector ($`\theta _{12},\mathrm{\Delta }m_{12}^2`$) or in the atmospheric sector ($`\theta _{23},\mathrm{\Delta }m_{23}^2`$). When studying the impact of solar parameter uncertainties we have fixed the atmospheric parameters to $`\theta _{23}=40^{}`$ and $`\mathrm{\Delta }m_{23}^2=2.5\times 10^3`$ eV<sup>2</sup> and computed four different grids, for $`s_{atm}=\pm 1;s_{oct}=\pm 1`$. When studying the impact of atmospheric parameter uncertainties we have fixed the solar parameters to $`\theta _{12}=32^{}`$ and $`\mathrm{\Delta }m_{12}^2=8.2\times 10^5`$ eV<sup>2</sup>. In this case, only two grids must be computed, one for each value of $`s_{atm}`$: the octant-degeneracy need not to be considered as an external (discrete) input, since $`\theta _{23}`$ is one of the free parameters in the grid.
When a three-parameters fit is performed, the other parameters in the grid are fixed to the corresponding present best fit value for $`\theta _{12},\mathrm{\Delta }m_{12}^2`$ and $`\mathrm{\Delta }m_{23}^2`$ or to $`\theta _{23}=40^{}`$ for the atmospheric angle (to take into account possible octant and mixed degeneracies, that would disappear for maximal mixing). This procedure is used to study the effect of one parameter at a time on the ($`\theta _{13},\delta `$) measure. The $`\chi ^2`$ function is:
$`\left[\chi ^2(\theta _{13},\delta ,x)\right]_{\alpha \beta }={\displaystyle \underset{\pm }{}}\left[{\displaystyle \frac{N_{\alpha \beta }^\pm (\theta _{13},\delta ,x;s_{atm},s_{oct})N_{\alpha \beta }^\pm (\overline{\theta }_{13},\overline{\delta },\overline{x};\overline{s}_{atm},\overline{s}_{oct})}{\delta N_{\alpha \beta }^\pm }}\right]^2,`$ (1)
with $`x`$ any of the parameters to be fitted in addition to $`\theta _{13}`$ and $`\delta `$, $`\pm `$ refers to neutrinos or antineutrinos and $`N_{\alpha \beta }^\pm `$ is the number of charged leptons $`l_\beta ^\pm `$ observed in the detector for a $`\nu _\alpha (\overline{\nu }_\alpha )`$ beam. The error on the sample $`N_{\alpha \beta }^\pm `$ is:
$$(\delta N_{\alpha \beta }^\pm )^2=\sigma _{N_{\alpha \beta }^\pm }^2+(ฯต_\beta ^\pm N_{\alpha \beta }^\pm )^2+(ฯต_\beta ^\pm B_{\alpha \beta }^\pm )^2,$$
where $`\sigma _{N_{\alpha \beta }^\pm }`$ is the statistical error on $`N_{\alpha \beta }^\pm `$ (Gaussian or Poissonian, depending on the corresponding statistics), $`B_{\alpha \beta }^\pm `$ is the sum of beam and detector backgrounds for the considered channel, computed as in Refs. , and $`ฯต_\beta ^\pm `$ is the total systematic error for the considered channel at a given facility. No covariance matrix for the non-fitted parameters has been considered. The three-parameters $`\chi ^2`$ function defines a three-dimensional 90% CL contour that is eventually projected onto the ($`\theta _{13},\delta `$) plane to perform a direct comparison with the standard two-parameters 90% CL contours for the considered setups<sup>4</sup><sup>4</sup>4A preliminary result obtained by means of this procedure has been presented in . .
A discussion on the statistical approach chosen and its difference with existing approaches is mandatory and can be found in App. A.
## 5 Parameter correlations and degeneracies
Parameter correlations and degeneracies arise in the determination of $`\theta _{13}`$ and $`\delta `$ at future neutrino experiments, as it has been studied in many papers -. The problem is due to the strong correlation between these two parameters in the appearance transition probabilities ($`\nu _e\nu _\mu ,\nu _\tau `$ and $`\nu _\mu \nu _e`$) and in the present (and near future) ignorance of two discrete unknowns, the sign of the atmospheric mass difference $`\mathrm{\Delta }m_{23}^2`$ and the $`\theta _{23}`$-octant, that can be parametrized by the sign variables $`s_{atm}=\text{sign}[\mathrm{\Delta }m_{23}^2]`$ and $`s_{oct}=\text{sign}[\mathrm{tan}(2\theta _{23})]`$ that take the values $`\pm 1`$ for $`\mathrm{\Delta }m_{23}^2>0(<0)`$ and $`\theta _{23}<45^{}(>45^{})`$, respectively. Solving the systems of equations corresponding to the four distinct choices of $`s_{atm}`$ and $`s_{oct}`$:
$`N_{\alpha \beta }^\pm (\overline{\theta }_{13},\overline{\delta };\overline{s}_{atm},\overline{s}_{oct})`$ $`=`$ $`N_{\alpha \beta }^\pm (\theta _{13},\delta ;s_{atm}=\overline{s}_{atm};s_{oct}=\overline{s}_{oct}),`$ (2)
$`N_{\alpha \beta }^\pm (\overline{\theta }_{13},\overline{\delta };\overline{s}_{atm},\overline{s}_{oct})`$ $`=`$ $`N_{\alpha \beta }^\pm (\theta _{13},\delta ;s_{atm}=\overline{s}_{atm},s_{oct}=\overline{s}_{oct}),`$ (3)
$`N_{\alpha \beta }^\pm (\overline{\theta }_{13},\overline{\delta };\overline{s}_{atm},\overline{s}_{oct})`$ $`=`$ $`N_{\alpha \beta }^\pm (\theta _{13},\delta ;s_{atm}=\overline{s}_{atm},s_{oct}=\overline{s}_{oct}),`$ (4)
$`N_{\alpha \beta }^\pm (\overline{\theta }_{13},\overline{\delta };\overline{s}_{atm},\overline{s}_{oct})`$ $`=`$ $`N_{\alpha \beta }^\pm (\theta _{13},\delta ;s_{atm}=\overline{s}_{atm},s_{oct}=\overline{s}_{oct}),`$ (5)
(with $`N_{\alpha \beta }^\pm `$ defined in the previous section) will result, in general, in the input pair ($`\overline{\theta }_{13},\overline{\delta }`$) plus seven additional solutions (the clones) to form an eightfold degeneracy: the intrinsic clone (Eq. 2), the sign clones (Eq. 3), the octant clones (Eq. 4) and the mixed clones (Eq. 5). A complete theoretical analysis of the clones location has been presented in Ref. , where an algorithm to numerically find each clone location in the ($`\theta _{13},\delta `$) plane as a function of the considered experimental setup and of the input parameters has been given. A similar approach can be applied to study the presence of degeneracies in the disappearance channels $`\nu _e\nu _e`$ and $`\nu _\mu \nu _\mu `$.
### 5.1 Correlation and degeneracies in $`\nu _e`$ disappearance
The $`\nu _e`$ disappearance probability does not depend on the CP violating phase $`\delta `$ and on the atmospheric $`\theta _{23}`$ mixing angle. The $`\theta _{13}`$ measurement is, therefore, not affected by $`(\theta _{13}\delta )`$ correlations nor by the $`s_{oct}`$ ambiguity. The $`\nu _e\nu _e`$ matter oscillation probability, expanded at second order in the small parameters $`\theta _{13}`$ and $`(\mathrm{\Delta }m_{12}^2L/E)`$ reads:
$`P_{ee}^\pm `$ $`=`$ $`1\left({\displaystyle \frac{\mathrm{\Delta }_{23}}{B_{}}}\right)^2\mathrm{sin}^2(2\theta _{13})\mathrm{sin}^2\left({\displaystyle \frac{B_{}L}{2}}\right)\left({\displaystyle \frac{\mathrm{\Delta }_{12}}{A}}\right)^2\mathrm{sin}^2(2\theta _{12})\mathrm{sin}^2\left({\displaystyle \frac{AL}{2}}\right),`$
where $`\mathrm{\Delta }_{23}=\mathrm{\Delta }m_{23}^2/2E`$, $`\mathrm{\Delta }_{12}=\mathrm{\Delta }m_{12}^2/2E`$, $`A=\sqrt{2}G_FN_e`$ and $`B_{}=|A\mathrm{\Delta }_{23}|`$ with $`\pm `$ for neutrinos (antineutrinos), respectively. This formula describes reasonably well the behaviour of the transition probability in the energy range covered by the considered $`\beta `$-Beam setup ($`L100`$ km and $`E_\nu `$ 100 MeV) and it illustrates clearly that two sources of ambiguities are still present in $`\nu _e`$ disappearance, $`s_{atm}`$ (for large values of $`\theta _{13}`$, i.e. in the โatmosphericโ region) and the $`(\theta _{13}\theta _{12})`$ correlation (for small values of $`\theta _{13}`$, i.e. in the โsolarโ region). A $`\beta `$-Beam could in principle improve our present errors on the solar parameters through $`\nu _e`$ disappearance. We have checked that this is not the case for the considered setup: at large $`\theta _{13}`$ the second term in eq. (LABEL:eq:disnue) dominates over the last term, that is more sensitive to solar parameters. On the other hand, for small $`\theta _{13}`$ the statistics is too low to improve present uncertainties on $`\theta _{12}`$ and $`\mathrm{\Delta }m_{12}^2`$ (remind that energy and baseline of the low-$`\gamma `$ $`\beta `$-Beam has not been chosen to fulfill this task, and therefore our results are not surprising at all). Eventually, in Ref. it has been shown that if systematic errors cannot be controlled better than at 5%, the $`\beta `$-Beam disappearance channel does not improve the CHOOZ bound on $`\theta _{13}`$.
Eq. (LABEL:eq:disnue) can be also applied to reactor experiments aiming to a precise measurement of $`\theta _{13}`$ in a โdegeneracy-freeโ environment. For the typical baseline and energy of a reactor experiment (e.g., $`L=1.05`$ km and $`E_\nu =4`$ MeV for the Double-Chooz proposal, ) we can safely consider antineutrino propagation in vacuum. As a consequence, no sensitivity to $`s_{atm}`$ is expected at these experiments, since $`B_{}\mathrm{\Delta }_{23}`$ for $`\mathrm{\Delta }_{23}A`$. It is very difficult that reactor experiments could test small values of $`\theta _{13}`$, and thus the $`\theta _{13}\theta _{12}`$ correlation (significant only in the โsolarโ region) can also be neglected.
### 5.2 Correlation and degeneracies in $`\nu _\mu `$ disappearance
A Super-Beam facility can perform an independent measurement of the atmospheric parameters via the $`\nu _\mu `$ disappearance channel: these kind of facilities should in principle reduce the error on the atmospheric mass difference to less than 10 % and on the atmospheric angle to $`10`$ % . It is thus interesting to study, as for the $`\nu _e`$ disappearance channel, the presence of parameter correlations and degeneracies. The vacuum oscillation probability expanded to the second order in the small parameters $`\theta _{13}`$ and $`(\mathrm{\Delta }_{12}L/E)`$ is:
$`P(\nu _\mu \nu _\mu )`$ $`=`$ $`1\left[\mathrm{sin}^22\theta _{23}s_{23}^2\mathrm{sin}^22\theta _{13}\mathrm{cos}2\theta _{23}\right]\mathrm{sin}^2\left({\displaystyle \frac{\mathrm{\Delta }_{23}L}{2}}\right)`$ (7)
$``$ $`\left({\displaystyle \frac{\mathrm{\Delta }_{12}L}{2}}\right)[s_{12}^2\mathrm{sin}^22\theta _{23}+\stackrel{~}{J}s_{23}^2\mathrm{cos}\delta ]\mathrm{sin}(\mathrm{\Delta }_{23}L)`$ (8)
$``$ $`\left({\displaystyle \frac{\mathrm{\Delta }_{12}L}{2}}\right)^2[c_{23}^4\mathrm{sin}^22\theta _{12}+s_{12}^2\mathrm{sin}^22\theta _{23}\mathrm{cos}(\mathrm{\Delta }_{23}L)],`$ (9)
where $`\stackrel{~}{J}=\mathrm{cos}\theta _{13}\mathrm{sin}2\theta _{12}\mathrm{sin}2\theta _{13}\mathrm{sin}2\theta _{23}`$. The first term in the first parenthesis is the dominant one and is symmetric under $`\theta _{23}\pi /2\theta _{23}`$. This is indeed the source of our present ignorance on $`s_{oct}`$. This symmetry is lifted by the other terms, that introduce a mild CP-conserving $`\delta `$-dependence also, albeit through subleading effects very difficult to isolate. We present our results for the $`\nu _\mu `$ disappearance channel in the ($`\theta _{23},\mathrm{\Delta }m_{23}^2`$) plane: as a consequence, we do not need to specify the $`\theta _{23}`$-octant, since the interval $`\theta _{23}[36^{},55^{}]`$ is spanned explicitly.
Solving the two systems of equations:
$$N_{\mu \mu }^\pm (\overline{\theta }_{23},\mathrm{\Delta }m_{atm}^2;\overline{s}_{atm})=N_{\mu \mu }^\pm (\theta _{23},|\mathrm{\Delta }m_{23}^2|;\overline{s}_{atm}),$$
(10)
$$N_{\mu \mu }^\pm (\overline{\theta }_{23},\mathrm{\Delta }m_{atm}^2;\overline{s}_{atm})=N_{\mu \mu }^\pm (\theta _{23},|\mathrm{\Delta }m_{23}^2|;\overline{s}_{atm}),$$
(11)
four different solutions are found for $`\overline{\theta }_{23}45^{}`$: two solutions from eq. (10), the input value $`\theta _{23}=\overline{\theta }_{23}`$ and $`\theta _{23}\pi /2\overline{\theta }_{23}`$, being the second solution not exactly at $`\theta _{23}=\pi /2\overline{\theta }_{23}`$ due to the small $`\theta _{23}`$-octant asymmetry; and two more solutions from eq. (11) at a different value of $`|\mathrm{\Delta }m_{23}^2|`$ . In eq. (9) we can see that changing sign to $`\mathrm{\Delta }m_{23}^2`$ the second term becomes positive: a change that must be compensated with an increase in $`|\mathrm{\Delta }m_{23}^2|`$ to give $`P_{\mu \mu }^\pm (\mathrm{\Delta }m_{atm}^2;\overline{s}_{atm})=P_{\mu \mu }^\pm (|\mathrm{\Delta }m_{23}^2|;\overline{s}_{atm})`$. The two solutions of eq. (11) corresponding to the wrong choice of $`s_{atm}`$ can be observed in Fig. 3(left), where equal-number-of-events (ENE) curves are computed for both $`s_{atm}=\overline{s}_{atm}`$ (solid) and $`s_{atm}=\overline{s}_{atm}`$ (dotted) at the considered Super-Beam facility. The two intersections are notably off the input pair $`\overline{\theta }_{23}=40^{}`$, $`\mathrm{\Delta }m_{23}^2=\mathrm{\Delta }m_{atm}^2`$. As expected, the two sign clones are located at $`|\mathrm{\Delta }m_{23}^2|\mathrm{\Delta }m_{atm}^2`$ and are almost symmetric with respect to $`\theta _{23}=45^{}`$. The shift in the vertical axis is a function of $`\theta _{13}`$ and $`\delta `$. If $`\overline{\theta }_{23}=45^{}`$ only two solutions (corresponding to different choices of $`s_{atm}`$) are expected.
We must also stress that such an uncertainty can be enhanced once we take into account that $`\theta _{13}`$ and $`\delta `$ are completely unknown (although the impact of this last parameter in $`\nu _\mu `$ disappearance is expected to be rather small). To gain some feeling on the precision that can be expected in a $`\nu _\mu `$ disappearance measurement at the SPL Super-Beam facility, we performed a full three-parameters analysis in $`\theta _{23},|\mathrm{\Delta }m_{23}^2|`$ and $`\theta _{13}`$ for the input parameters $`\overline{\theta }_{23}=40^{}`$, $`\mathrm{\Delta }m_{23}^2=2.5\times 10^3`$ eV<sup>2</sup> and $`\overline{\theta }_{13}=7^{}`$. We have then projected the 90% three-parameters CL contour onto the $`(\theta _{23},|\mathrm{\Delta }m_{23}^2|)`$ plane, Fig. 3(right). In the absence of a complete simulation of the systematics and the background for the $`\nu _\mu `$ disappearance channel at the SPL Super-Beam , we have adopted as an estimate of the expected background and efficiency those used in and for the $`\nu _e\nu _\mu `$ appearance channel at the $`\beta `$-Beam facility. A 2% systematic error has been assumed. The solid line refers to the projection of the three-dimensional 90 % CL contour on the ($`\theta _{23},|\mathrm{\Delta }m_{23}^2|`$) plane for $`s_{atm}=\overline{s}_{atm}`$, the dotted line to the projection of the 90 % CL contour for $`s_{atm}=\overline{s}_{atm}`$.
As expected, the three-parameters fit presents a second allowed region in the parameter space at $`|\mathrm{\Delta }m_{23}^2|>\mathrm{\Delta }m_{atm}^2`$ when the wrong $`s_{atm}`$ is considered. Notice that, performing a three-parameters fit in $`\theta _{23},\mathrm{\Delta }m_{23}^2`$ and $`\delta `$, the difference between the two- and three-parameters contours is much smaller. In it has been shown that a larger spread in $`\theta _{23}`$ is found for $`\overline{\theta }_{23}45^{}`$. We can perform fits with $`\theta _{23}`$ non-maximal and for different input values for $`\overline{\theta }_{13}[0,10^{}]`$. The result of such an analysis is that the SPL Super-Beam will be able to measure $`\theta _{23}`$ in the interval $`[36^{},55^{}]`$ and $`\mathrm{\Delta }m_{23}^2`$ in $`[2.3,2.9]\times 10^3\mathrm{eV}^2`$ for the input pair $`\overline{\theta }_{23}=40^{},\mathrm{\Delta }m_{23}^2=\mathrm{\Delta }m_{atm}^2`$. Notice that the expected SPL precision on $`\mathrm{\Delta }m_{23}^2`$ is comparable with what expected at T2K-I . On the other hand, the expected SPL precision on $`\theta _{23}`$ is much worse than the T2K-I one, a consequence of the fact that the considered SPL setup is a counting experiment and it has no energy resolution.
## 6 Impact of parameter uncertainties on $`\theta _{13}`$ and $`\delta `$
In this section, we discuss the impact of the uncertainties in the solar and atmospheric parameters to the simultaneous measurement of $`\theta _{13}`$ and $`\delta `$ at the considered $`\beta `$-Beam, Super-Beam and Neutrino Factory. The three experiments will be discussed separately.
In all fits we have combined informations from all available channels for both polarities,
$$\chi ^2(\theta _{13},\delta ,x)=\underset{i}{}\chi _i^2(\theta _{13},\delta ,x),$$
(12)
where $`\chi _i^2`$ is the three-parameters $`\chi ^2`$ function defined as in eq. (1) for a given channel and polarity. All channels have been taken as independent measurement and no covariance matrix has been introduced, following the approach described in Sect. 4. In general, a โpessimisticโ systematic error, $`ฯต^\pm =5`$%, has been used in appearance channels. On the other hand, a 2% systematic error has been used in disappearance channels.
### 6.1 The solar sector
We study the effect of present uncertainties on the solar sector parameters in the measurement of $`\theta _{13}`$ and $`\delta `$ performing two distinct three-parameters fit in $`\theta _{13},\delta `$ and $`\theta _{12}`$ (for fixed $`\mathrm{\Delta }m_{12}^2`$) or $`\mathrm{\Delta }m_{12}^2`$ (for fixed $`\theta _{12}`$). The fits have been performed using 10 years of $`\beta `$-Beam running with both polarities.
The projection of the three-parameters 90% CL contours onto the ($`\theta _{13},\delta `$) plane are presented in Fig. 4. In the left panel we have fixed $`\mathrm{\Delta }m_{12}^2=8.2\times 10^5`$ eV<sup>2</sup> and drawn the projection of the three-dimensional contours for $`\chi ^2(\theta _{13},\delta ,\theta _{12})`$, for $`\theta _{12}[29^{},36^{}]`$. In the right panel we have fixed $`\theta _{12}=32^{}`$ and drawn the projection of the three-dimensional contours for $`\chi ^2(\theta _{13},\delta ,\mathrm{\Delta }m_{12}^2)`$, for $`\mathrm{\Delta }m_{12}^2[7.5,9.1]\times 10^5`$ eV<sup>2</sup>. In both cases, the atmospheric parameters have been fixed to $`\mathrm{\Delta }m_{23}^2=2.5\times 10^3`$ eV<sup>2</sup> and $`\theta _{23}=40^{}`$. The input values for the two unknowns are $`\overline{\theta }_{13}=2^{},7^{}`$ and $`\overline{\delta }=45^{}`$. For each panel, the results for the four different choices of the two discrete variables, $`s_{atm}`$ and $`s_{oct}`$, are presented separately. Finally, the projection of the three-parameters 90% CL contours (solid lines) are directly compared with the two-parameters 90% CL contours (dashed lines) obtained fixing the solar parameters to their present best fit values, $`\theta _{12}=32^{},\mathrm{\Delta }m_{12}^2=8.2\times 10^5`$ eV<sup>2</sup>.
As we can see in both panels, most of the plotted three-parameters contours coincide for any practical purpose with the corresponding two-parameters ones, with small deviations easily explained by the different CL in two- and three-parameters $`\chi ^2`$. As a result, we claim that the impact of solar parameter uncertainties on the measurement of $`\theta _{13}`$ and $`\delta `$ is negligible for $`\overline{\theta }_{13}2^{}`$. This is indeed a consequence of the subleading dependence of the $`\nu _e\nu _\mu `$ oscillation probability on the solar parameters (see, for example, Refs. ) for large values of $`\overline{\theta }_{13}`$. When $`\overline{\theta }_{13}`$ is large, we are in what has been called the โatmospheric regimeโ in Ref. ; only for $`\overline{\theta }_{13}`$ below the verge of the $`\beta `$-Beam $`\theta _{13}`$-sensitivity, i.e. for $`\overline{\theta }_{13}>2^{}`$, we enter in the so-called โsolar regimeโ. Clearly, no signal is expected at the $`\beta `$-Beam in this case: we can thus safely claim that the solar parameter uncertainties do not affect significantly the measurement of $`\theta _{13}`$ and $`\delta `$ at the considered facility.
Similar conclusions can be drawn for the SPL Super-Beam and for different values of $`\overline{\delta }`$ and will therefore not be repeated here. For the rest of the paper, the solar parameters will be considered as fixed external inputs: $`\theta _{12}=32^{}`$ and $`\mathrm{\Delta }m_{12}^2=8.2\times 10^5`$ eV<sup>2</sup>.
### 6.2 The atmospheric sector at the $`\beta `$-Beam
As for the solar sector, we study the effect of present uncertainties on the atmospheric sector parameters in the measurement of $`\theta _{13}`$ and $`\delta `$ performing two distinct three-parameters fit in $`\theta _{13},\delta `$ and $`\theta _{23}`$ (for fixed $`\mathrm{\Delta }m_{23}^2`$) or $`\mathrm{\Delta }m_{23}^2`$ (for fixed $`\theta _{23}`$).
The comparison between two- and three-parameters fits is presented in Fig. 5, where the projection of the three-parameters 90% CL contours onto the ($`\theta _{13},\delta `$) plane (solid lines) and the two-parameters 90% CL contours (dashed lines) have been plotted separately for each possible choice of the two discrete variables, $`s_{atm}`$ and $`s_{oct}`$. In the left panel we have fixed $`\mathrm{\Delta }m_{atm}^2=2.5\times 10^3`$ eV<sup>2</sup> and drawn the projection of the three-dimensional contours for $`\chi ^2(\theta _{13},\delta ,\theta _{23})`$, for $`\theta _{23}[36^{},55^{}]`$ (see Tab. 1). In the right panel we have fixed $`\theta _{23}=40^{}`$ and drawn the projection of the three-dimensional contours for $`\chi ^2(\theta _{13},\delta ,\mathrm{\Delta }m_{23}^2)`$, for $`\mathrm{\Delta }m_{23}^2[1.7,3.5]\times 10^3`$ eV<sup>2</sup> (see Tab. 1). The input values for the two unknowns are $`\overline{\theta }_{13}=2^{},7^{}`$ and $`\overline{\delta }=45^{}`$. The two-parameters contours have been drawn using fixed values for the atmospheric parameters, $`\theta _{23}=40^{},\mathrm{\Delta }m_{atm}^2=2.5\times 10^3`$ eV<sup>2</sup>. We have checked that no significant improvement is observed when the systematic error in the appearance channel is set to 2%.
In Fig. 5 it is manifest the impact of atmospheric parameter uncertainties on the measurement of $`\theta _{13}`$ and $`\delta `$, for both $`\theta _{23}`$ and $`\mathrm{\Delta }m_{23}^2`$ fits. Both unknowns are measured with errors much larger than those expected from two-parameters contours. This must be compared with the results of the previous section, where it has been shown that solar parameter uncertainties have a negligible impact.
Consider first the left panel of Fig. 5: the results from a three-parameters fit in $`(\theta _{13},\delta ,\theta _{23})`$. Notice that, being $`\theta _{23}`$ a fitting variable in the whole range $`\theta _{23}[36^{},55^{}]`$, the effect of the โoctant ambiguityโ is automatically taken into account by the three-parameters $`\chi ^2`$ function. For this reason the three-parameters contours labelled as โtrueโ and โoctantโ are identical. Notice also that this is not the case for the two-parameters contours, where the choice of $`s_{oct}`$ is reflected in contours located at different values of $`\theta _{13}`$ with respect to the input $`\overline{\theta }_{13}`$. For $`\overline{\theta }_{13}=7^{}`$ we can see that a large error in $`\theta _{13}`$ is induced by the uncertainty in $`\theta _{23}`$ (with $`\mathrm{\Delta }\theta _{13}`$ as large as $`4^{}`$). This is a consequence of the fact that the leading term in the $`\nu _e\nu _\mu `$ oscillation probability is proportional to the combination $`\mathrm{sin}^2(2\theta _{13})\mathrm{sin}^2\theta _{23}`$: to compensate a change in $`\theta _{23}`$, a change in $`\theta _{13}`$ is needed. For smaller values of $`\overline{\theta }_{13}`$ this effect is much smaller. The spread in $`\delta `$ is, on the other hand, extremely similar in two- and three-parameters contours. This is a consequence of the fact that a $`\delta `$-dependence would be induced in the fit through the subleading term in the oscillation probability, that is proportional to $`\mathrm{sin}(2\theta _{23})`$ and thus less sensitive to changes in $`\theta _{23}`$ in the almost symmetric interval considered. Notice that three-parameters contours have a box-like shape, with no strong $`\theta _{13}\delta `$ correlation. Finally, the largest values of $`\theta _{13}`$ are observed for both choices of $`s_{atm}`$ at lower values of $`\theta _{23}`$, whereas the smallest values of $`\theta _{13}`$ are reached for larger values of $`\theta _{23}`$.
Consider now the right panel of Fig. 5: the results from a three-parameters fit in $`(\theta _{13},\delta ,\mathrm{\Delta }m_{23}^2)`$. In this case to different choices of $`s_{oct}`$ correspond different contours ($`\theta _{23}`$ is a fixed external input and not a free parameter in the fit). As in the previous case, a large error in $`\theta _{13}`$ is induced by the error on the atmospheric parameter, especially for $`\overline{\theta }_{13}=7^{}`$. The largest value of $`\theta _{13}`$ is associated to the smallest value of $`\mathrm{\Delta }m_{23}^2`$, in all plots. A characteristic feature of this fit is the significant $`\delta `$-dependence that can be observed in all plots and that was not present in the fits in $`\theta _{23}`$. The three-parameters 90% CL contours have a triangular shape (the error in $`\delta `$ reduces for large values of $`\theta _{13}`$), pointing to a strong $`\theta _{13}\delta `$ correlation. It is important to stress that, when $`\mathrm{\Delta }m_{23}^2`$ is the free parameter, the overall error in $`\delta `$ is significantly larger in the three- than in the two-parameters contours. For both values of $`\overline{\theta }_{13}=2^{},7^{}`$ roughly half of the $`\delta `$-parameter space is covered. This result, considerably worse than what expected from two-parameters fit, can be compared with the CP-coverage expectation (explained in App. A) for this particular input pair, Fig. 11. For both methods, a rather large error in $`\delta `$ is indeed expected.
It is clear from Fig. 5 that both atmospheric parameter uncertainties are extremely important in the measurement of $`\theta _{13}`$ and $`\delta `$: the three-parameters 90% CL allowed regions are considerably worse than those obtained with two-parameters fits. In particular, for the shown input pairs, $`\delta `$ would remain completely unknown in the interval $`\delta [0,\pi ]`$ and $`\theta _{13}`$ would be known only with a large error.
As an example of how the situation can be improved when using reduced uncertainties on the atmospheric parameters, in Fig. 13 we present the projection of the three-dimensional 90% CL contours onto the ($`\theta _{13},\delta `$) plane using the expected uncertainties on the atmospheric parameters after T2K-I (last column of Tab. 1): $`\theta _{23}[38^{},43^{}][48^{},52^{}]`$ and $`\mathrm{\Delta }m_{23}^2[2.42,2.61]\times 10^3`$ eV<sup>2</sup> for $`s_{atm}=+`$ and $`\mathrm{\Delta }m_{23}^2[2.46,2.64]\times 10^3`$ eV<sup>2</sup> for $`s_{atm}=`$, . The octant-ambiguity (that will not be solved at T2K-I) recover its discrete nature: separate regions of the parameter space will be spanned by different choices of $`s_{oct}`$. In this case, all choices of the two discrete variables $`s_{atm}`$ and $`s_{oct}`$ are presented together and no comparison with two-parameters contours is shown. In top panels $`x=\theta _{23}`$; in bottom panels $`x=\mathrm{\Delta }m_{23}^2`$. The results of the three-parameters fit with expected uncertainties (right panels), are directly compared with the results presented in Fig. 5 computed with the present uncertainties (left panels). The reduction of the uncertainties on the atmospheric parameters has indeed an important effect on the measurement of $`\theta _{13}`$ and $`\delta `$. As it can be seen in the right panels of Fig. 13 a significant reduction of the $`\theta _{13}`$-spread is achieved, with plots resembling those obtained with standard two-parameters contours and fixed external atmospheric parameters (see Refs. ). The $`\delta `$-spread is also reduced considerably with respect to the results obtained with present uncertainties. These comments apply to both $`\overline{\theta }_{13}=2^{},7^{}`$. Notice that, as expected being $`\theta _{23}`$ restricted to one octant only, the octant- and mixed- ambiguities show themselves as separate contours in the ($`\theta _{13},\delta `$) plane, as for two-parameters fits.
A final comment on the impact of the uncertainties on the atmospheric parameters on the measurement of $`\theta _{13}`$ and $`\delta `$ at the low-gamma $`\beta `$-Beam is in order. We have shown that, with present uncertainties, the measurement of the two unknowns in the PMNS mixing matrix is severely spoiled. Errors as large as $`\mathrm{\Delta }\theta _{13}4^{}`$ are found, and half of the parameter space in $`\delta `$ is spanned for different values of $`\overline{\delta }`$. This corresponds to a CP-coverage $`\xi 0.5`$, a value that spoils completely the possibility to distinguish a CP-violating signal from a CP-conserving one at the considered facility (see App. A). A significant reduction in the uncertainties on the atmospheric parameters is mandatory if we plan to use such a facility to look for $`\delta `$. If $`\theta _{23}`$ and $`\mathrm{\Delta }m_{23}^2`$ can be measured at the T2K-I experiment with the expected precision and for any value of $`\overline{\theta }_{23}`$, only then the results of present two-parameters studies for facilities of this kind can be considered reliable.
In App. B we present the results for different choices of $`\overline{\delta }`$, to illustrate the generality of the results above.
### 6.3 The atmospheric sector at the Super-Beam
We now repeat the analysis of the impact of present and expected uncertainties on the atmospheric sector parameters in the measurement of $`\theta _{13}`$ and $`\delta `$ at a different facility: the SPL Super-Beam. Again, two distinct three-parameters fits in $`\theta _{13},\delta `$ and $`\theta _{23}`$ (for fixed $`\mathrm{\Delta }m_{23}^2`$) or $`\mathrm{\Delta }m_{23}^2`$ (for fixed $`\theta _{23}`$) have been performed, with the Super-Beam running 2 years with $`\pi ^+`$ and 8 years with $`\pi ^{}`$ to accumulate comparable statistics for neutrinos and antineutrinos. A significant difference between this facility and the low-gamma $`\beta `$-Beam considered previously is in that the $`\nu _\mu `$ disappearance channel at the Super-Beam reduces the uncertainties on the atmospheric parameters, as it can be seen in Fig. 3 (whereas the $`\nu _e`$ disappearance channel is useless to this purpose, see Ref. ). In this case we therefore do not present results using โpresentโ and โexpectedโ uncertainties, but we just combine the results from the appearance and disappearance channel. We have checked that using a โpessimisticโ 5% systematic error in the disappearance channel does not change significantly our results.
The comparison between two- and three-parameters fits is presented in Fig. 6. In the left panel we have fixed $`\mathrm{\Delta }m_{atm}^2=2.5\times 10^3`$ eV<sup>2</sup> and drawn the projection of the three-dimensional contours for $`\chi ^2(\theta _{13},\delta ,\theta _{23})`$, for $`\theta _{23}[36^{},55^{}]`$ (see Tab. 1). In the right panel we have fixed $`\theta _{23}=40^{}`$ and drawn the projection of the three-dimensional contours for $`\chi ^2(\theta _{13},\delta ,\mathrm{\Delta }m_{23}^2)`$, for $`\mathrm{\Delta }m_{23}^2[2.3,2.9]\times 10^3`$ eV<sup>2</sup> (see Tab. 1 and Sect. 5.2). The input values for the two unknowns are $`\overline{\theta }_{13}=2^{},7^{}`$ and $`\overline{\delta }=45^{}`$.
The main difference between two- and three-parameters contours resides in that in the latter we observe some clones absent in the two-parameters plots. This is a consequence of the not satisfactory expected improvement on the error in $`\theta _{23}`$ and $`\mathrm{\Delta }m_{23}^2`$ for $`\theta _{23}=40^{}`$.
Also in this case the measurement of $`\theta _{13}`$ and $`\delta `$ is severely affected by the uncertainties on the atmospheric parameters. Somewhat smaller errors are found in $`\theta _{13}`$ and $`\delta `$ with respect to the $`\beta `$-Beam case, but still almost half of the parameter space in $`\delta `$ is spanned for different values of $`\overline{\delta }`$. A crucial point is that it does not seem that the $`\nu _\mu `$ disappearance channel is capable of a significant reduction in the error on the atmospheric mixing angle $`\theta _{23}`$. The T2K-I experiment will therefore be crucial, if indeed the expected precision in the atmospheric angle can be met for any value of $`\theta _{23}`$.
In App. B we present the results for different choices of $`\overline{\delta }`$, to illustrate the generality of the results above.
### 6.4 The atmospheric sector at the Neutrino Factory
We repeat the analysis of the impact of atmospheric parameters uncertainties in the measurement of $`\theta _{13}`$ and $`\delta `$ at a third facility: the CERN-based SPL-fuelled 50 GeV Neutrino Factory. We want to show in this way how the results of Sects. 6.2 and 6.3 are quite general and must be taken into account at any facility that is considered when looking for $`\theta _{13}`$ and $`\delta `$.
As before, two distinct three-parameters fit in $`\theta _{13},\delta `$ and $`\theta _{23}`$ (for fixed $`\mathrm{\Delta }m_{23}^2`$) or $`\mathrm{\Delta }m_{23}^2`$ (for fixed $`\theta _{23}`$) have been performed, with the Neutrino Factory running 5 years with $`\mu ^+`$ and 5 years with $`\mu ^{}`$.
In the absence of an updated analysis of the expected reduction of atmospheric parameters uncertainties at this facility through $`\nu _\mu \nu _\mu ,\nu _e\nu _e`$ and $`\nu _\mu \nu _\tau `$ (see Ref. for old analyses), we only present results combining the two appearance channels $`\nu _e\nu _\mu `$ (i.e. the โgoldenโ channel) and $`\nu _e\nu _\tau `$ (i.e. the โsilverโ channel) for both polarities. We use the expected uncertainties after T2K-I, in order to get a preliminar understanding of the impact of atmospheric parameter uncertainties at this facility.
The comparison between two- and three-parameters fits is presented in Fig. 7. In the left panel we have fixed $`\mathrm{\Delta }m_{atm}^2=2.5\times 10^3`$ eV<sup>2</sup> and drawn the projection of the three-dimensional contours for $`\chi ^2(\theta _{13},\delta ,\theta _{23})`$, for $`\theta _{23}[38^{},43^{}][48^{},52^{}]`$ (see Tab. 1). In the right panel we have fixed $`\theta _{23}=40^{}`$ and drawn the projection of the three-dimensional contours for $`\chi ^2(\theta _{13},\delta ,\mathrm{\Delta }m_{23}^2)`$, for $`\mathrm{\Delta }m_{23}^2[2.4,2.7]\times 10^3`$ eV<sup>2</sup> (see Tab. 1 and Sect. 5.2). The input values for the two unknowns are $`\overline{\theta }_{13}=2^{},7^{}`$ and $`\overline{\delta }=42^{}`$.
First of all notice that at the Neutrino Factory the sign and mixed degeneracies are solved, being the magnetized iron detector with a $`L=3000`$ km baseline extremely sensitive to matter effects and thus capable to measure $`s_{atm}`$. For this reason we only present two panels, corresponding to two possible choices of the $`\theta _{23}`$-octant, $`s_{oct}=\pm \overline{s}_{oct}`$. As for the SPL Super-Beam, for small $`\overline{\theta }_{13}`$ the two- and three-parameters contours practically coincide. On the other hand, for $`\overline{\theta }_{13}`$ large we must make a distinction between $`s_{oct}=\overline{s}_{oct}`$ and $`s_{oct}=\overline{s}_{oct}`$: whereas the impact of the atmospheric uncertainties for $`s_{oct}=\overline{s}_{oct}`$ is marginal (something already observed in , where the covariance matrix approach was adopted and only the right choice of the $`\theta _{23}`$-octant was considered), we notice how extra octant clones are present in the three-parameters contours that are absent in the two-parameters ones when the wrong choice of $`s_{oct}`$ is taken. This happens because in the three-dimensional parameter space $`\theta _{23}`$ cooperates with $`\theta _{13}`$ to identify a low $`\chi ^2`$ region with $`\delta \overline{\delta }`$ but with $`\theta _{13}<\overline{\theta }_{13}`$.
As for the other facilities, we have seen that the impact of the uncertainties on the atmospheric parameters on the measurement of $`\theta _{13}`$ and $`\delta `$ at the Neutrino Factory is relevant (albeit perhaps not as important as for the $`\beta `$-Beam and the Super-Beam previously discussed). Again, we stress that the loss in precision is more important for large $`\overline{\theta }_{13}`$ than for small $`\overline{\theta }_{13}`$, a region of the PMNS parameter space that will be selected or excluded by the approaching T2K-I experiment. This is indeed a crucial problem for precision measurements of the PMNS matrix elements.
In App. B we present the results for different choices of $`\overline{\delta }`$, to illustrate the generality of the results above.
## 7 CP-violation discovery potential
Eventually, in Figs. 8-10 we compare the sensitivity to $`(\theta _{13},\delta )`$ obtained with a two-parameters fit in ($`\theta _{13},\delta `$) or a three-parameters fit in ($`\theta _{13},\delta ,\theta _{23}`$) or ($`\theta _{13},\delta ,\mathrm{\Delta }m_{23}^2`$) at the three considered facilities. The 3$`\sigma `$ contours have been computed as in Ref. : at a fixed $`\overline{\theta }_{13}`$, we look for the smallest (largest) value of $`|\overline{\delta }|`$ for which the two- (three-) parameters 3$`\sigma `$ contours of any of the degenerate solutions (true, sign, octant and mixed) do not touch $`\delta =0^{}`$ nor $`\delta =180^{}`$. Notice that, although the input $`\overline{\theta }_{13}`$ value is fixed, the clones can touch $`\delta =0^{},180^{}`$ at $`\theta _{13}\overline{\theta }_{13}`$, also<sup>5</sup><sup>5</sup>5This is not the case of Fig. 11 in Ref. , where the excluded region in $`\delta `$ at fixed $`\overline{\theta }_{13}`$ in the absence of a CP-violating signal at 90% CL is presented. In practice, in that figure we compare $`N_\pm (\overline{\theta }_{13},\delta )`$ with $`N_\pm (\overline{\theta }_{13},0^{})`$, thus obtaining a one-parameter sensitivity plot in $`\delta `$ only.. The outcome of this procedure is finally plotted, representing the region in the $`(\theta _{13},\delta )`$ parameter space for which a CP-violating signal is observed at 3$`\sigma `$. Within this approach we can thus take fully into account the impact of the parameter degeneracies in the CP-violation discovery potential of the three facilities. As for the previous section, we have applied a 2% systematic error on disappearance channels and a 5% systematic error on appearance channels. As in the previous section we used the expected errors on the atmospheric parameters after T2K-I for the three-parameters fits at the $`\beta `$-Beam and the Neutrino Factory. The SPL Super-Beam analysis relies on SPL data, only (see Sect. 5.2).
Notice that results are given for the whole allowed range in $`\delta `$, $`\delta [180^{},180^{}]`$. This is particularly appropriate, since only an approximate symmetry is observed for $`|\delta |\pi /2`$ and $`|\delta |\pi /2`$ and no symmetry at all between positive and negative $`\delta `$ in the case of the $`\beta `$-Beam and of the Neutrino Factory.
Consider first Fig. 8, that refers to $`\beta `$-Beam results. Notice that the discovery potential is not symmetric for positive and negative values of $`\delta `$, as it has already been observed in Ref. . This asymmetric behaviour of the $`\beta `$-Beam is indeed a statistical mirage caused by the low background in the appearance antineutrino sample and the high background in the appearance neutrino one (see Ref. ). A proper statistical treatment should be performed, following Ref. , to get rid of this asymmetry for small $`\mathrm{sin}^2\theta _{13}`$: the treatment, however, is extremely time-consuming and we do not consider meaningful applying it here. A further asymmetry can be observed in the different behaviour of the three-parameters 3$`\sigma `$ contours projection onto the ($`\theta _{13},\delta `$) plane for positive and negative values of $`\delta `$: whereas for $`\delta >0`$ we observe that the smallest value of $`\mathrm{sin}^2\theta _{13}`$ for which a CP-violating phase can be distinguished from a null result goes from $`[\mathrm{sin}^2\theta _{13}]_{min}=3\times 10^45\times 10^4`$, for $`\delta <0`$ we get $`[\mathrm{sin}^2\theta _{13}]_{min}=2\times 10^34\times 10^3`$ for both the $`\theta _{23}`$ and $`|\mathrm{\Delta }m_{23}^2|`$ fits. A small loss in the discovery potential of this facility with respect to the two-parameters fit is observed in both three-parameters fits for negative $`\delta `$. In particular, the region in which a CP-violating signal can be distinguished from a CP-conserving one goes from $`\delta [80^{},130^{}]\delta [90^{},120^{}]`$.
Consider now Fig. 9, that refers to Super-Beam results. The strong asymmetry for positive and negative $`\delta `$ is not observed in this case, both for two- and three-parameters fits. The impact of the third fitting variable, being it $`\theta _{23}`$ or $`\mathrm{\Delta }m_{23}^2`$, is a rather small loss in the minimum value of $`\mathrm{sin}^2\theta _{13}`$ for which a CP-violating phase is distinguished from a null result: $`[\mathrm{sin}^2\theta _{13}]_{min}=9\times 10^41.2\times 10^3`$ for $`x=\theta _{23}`$ and $`[\mathrm{sin}^2\theta _{13}]_{min}=9\times 10^41.4\times 10^3`$ for $`x=\mathrm{\Delta }m_{23}^2`$. The loss in the $`\delta `$-interval that is distinguishable from a null result is rather small for three-parameters fits in $`\mathrm{\Delta }m_{23}^2`$ and negligible when fitting in $`\theta _{23}`$.
For the Neutrino Factory, Fig. 10, we observe a mixed situation: a strong asymmetry between positive and negative $`\delta `$ regions (as for the $`\beta `$-Beam), but a very small difference between two- and three-parameters fits (as for the Super-Beam). The asymmetry, however, it is not a consequence of asymmetric signal-to-noise ratios<sup>6</sup><sup>6</sup>6It must also be reminded that for 5 years of data taking in each polarity, a smaller statistics is accummulated in the wrong-sign muon sample for initial $`\mu ^{}`$ than for initial $`\mu ^+`$, due to the different $`\nu N`$ and $`\overline{\nu }N`$ cross-sections. This reduces the sensitivity to $`\theta _{13}`$ for negative $`\overline{\delta }`$. as for the $`\beta `$-Beam but, rather, of a โparametric conspiracyโ that for the chosen values of energy and baseline results in clones that for many negative values of $`\overline{\delta }`$ move toward $`\delta =0^{}`$ or $`\delta =180^{}`$ , thus preventing a clean identification of a CP-violating signal (see Figs. 22-26). The impact of the third fitting variable, being it $`\theta _{23}`$ or $`\mathrm{\Delta }m_{23}^2`$, is a rather small loss in the minimum value of $`\mathrm{sin}^2\theta _{13}`$ for which a CP-violating phase is distinguished from a null result for negative $`\overline{\delta }`$: $`[\mathrm{sin}^2\theta _{13}]_{min}=1.5\times 10^32\times 10^3`$. On the other hand, for positive values of $`\overline{\delta }`$, $`[\mathrm{sin}^2\theta _{13}]_{min}=2.5\times 10^4`$ both for two- and three-parameters fits.
## 8 Conclusions
The simultaneous measurement of $`\theta _{13}`$ and $`\delta `$ has been often performed in the literature considering the solar and atmospheric PMNS parameters as external quantities fixed to their best fit values. This is an approximation that has been adopted to get a first insight on the problems related to the ($`\theta _{13},\delta `$) measurement. The experimental uncertainties on these parameters can in principle affect the measurement of the unknowns, and it seemed important to us to perform an analysis that could go beyond the two-parameters fits presented in the literature.
In this paper we therefore have tried to study the impact that solar and atmospheric sector parameter uncertainties have on the measurement of $`\theta _{13}`$ and $`\delta `$ at three out of the many proposed setups, the standard low-$`\gamma `$ $`\beta `$-Beam, the 4 MWatt SPL Super-Beam and the 50 GeV SPL-fuelled Neutrino Factory. By doing this we wanted to catch the characteristic features of the inclusion of external parameters uncertainties in a ($`\theta _{13},\delta `$) measurement.
Our first goal has been to identify which of the external parameters affects the most the results of two-parameters fits. To do so we have performed a series of three-parameters fits in $`\theta _{13},\delta `$ and one of the other parameters ($`\theta _{12},\mathrm{\Delta }m_{12}^2,\theta _{23}`$ and $`\mathrm{\Delta }m_{atm}^2`$) in turn as the third fitting variable and compared our results with standard two-parameters fits. It turned out that the impact of solar parameters uncertainties on the measurement of ($`\theta _{13},\delta `$) is negligible, in practice, whereas present uncertainties on the atmospheric parameters are large enough to modify in a significant way the results of two-parameters fits. In particular, we have noticed that the main cause of the worsening from two- to three-parameters fits are the wide displacements of the so-called clones, parametric degeneracies due to multiple solutions of eqs. (2)-(5), as a consequence of small changes in the external parameters. These results are general to all the considered facilities.
We have then focused our attention on how the reduction of the atmospheric parameters uncertainties could ameliorate the previous results. To this respect, the three facilities we have considered are on different footing. On one side, the $`\nu _e`$ disappearance channel at the standard low-$`\gamma `$ $`\beta `$-Beam cannot improve on its own the present measurement of the atmospheric parameters. This facility, therefore, must rely on other experiments to meet its goal on $`\theta _{13}`$ and $`\delta `$. Luckily enough, it turns out that the precision on $`\theta _{23}`$ and $`\mathrm{\Delta }m_{23}^2`$ expected at the approved T2K-phase I experiment, if met, would be enough to improve our three-parameters fits and reproduce the results of two-parameters fits in the literature (that, however, were not so good). On the other hand, we have shown that the $`\nu _\mu `$ disappearance channel at the 4 MWatt SPL Super-Beam does improve the present errors on the atmospheric parameters. This facility, therefore, should not necessarily rely on external inputs. The combination of appearance and disappearance data, indeed, improve significantly our three-parameters fits. However, $`\theta _{23}`$ is not measured well enough and extra clones are still present in the ($`\theta _{13},\delta `$) that are absent in two-parameters contours. Finally, the Neutrino Factory has certainly the potential to improve significantly the precision on the atmospheric parameters through $`\nu _e`$ and $`\nu _\mu `$ disappearance and the $`\nu _\mu \nu _\tau `$ appearance channel, something that we have not studied in this paper. Using the errors on $`\theta _{23}`$ and $`\mathrm{\Delta }m_{23}^2`$ expected at the T2K-I experiment we have checked that extra clones are present in three-parameters fits that were absent in the two-parameters analysis. This is a clear indication of the fact that the problem we are addressing is common to all the facilities, not only to the low-$`\gamma `$ $`\beta `$-Beam or the SPL Super-Beam. It is not sufficient to just wait and see, but it must be taken into account when envisaging future facilities to look for $`\theta _{13}`$ and $`\delta `$.
To include the impact of external parameters uncertainties, other methods than direct multi-parameter fits have been proposed in the literature. For this reason, we have presented a direct comparison of our three-parameters fit results with the so-called CP-coverage introduced in Ref. . We have shown that in both methods a significant worsening of two-parameters fits arise as a consequence of the inclusion of errors on the external parameters in the fit. Whereas the CP-coverage method, however, can be quite useful to condense informations about the CP-sensitivity of a facility irrespectively of the specific input pair ($`\overline{\theta }_{13},\overline{\delta }`$) considered, direct three-parameters fits offer a detailed information for both $`\theta _{13}`$ and $`\delta `$ for specific points in the parameter space. We believe that the two methods are, in some sense, complementary and should be combined to get a thorough view of the performance of a specific facility designed to measure the ($`\theta _{13},\delta `$) pair. To this scope we have presented in App. B the results of a series of three-parameters fits for the standard low-$`\gamma `$ $`\beta `$-Beam, the 4 MWatt SPL Super-Beam and the 50 GeV Neutrino Factory for different choices of the input parameters.
Eventually, we have studied the impact of the atmospheric parameters uncertainties in the CP-violation discovery potential of the three considered facilities. Our results show that the discovery potential at the standard low-$`\gamma `$ $`\beta `$-Beam and at the SPL Super-Beam is somewhat reduced for negative values of $`\delta `$ when uncertainties on $`\theta _{23}`$ and $`\mathrm{\Delta }m_{23}^2`$ are taken into account. On the other hand the Neutrino Factory appears less affected by the inclusion of external parameter errors.
In conclusion, we think that this paper shows that present uncertainties on atmospheric parameters are indeed too large so that the widely adopted approximation of fixing $`\theta _{23}`$ and $`\mathrm{\Delta }m_{23}^2`$ to their present best fit values be reliable. A new phase of experiments that could improve these uncertainties are needed. The precision that is expected on the atmospheric parameters at the T2K-I experiment is shown to be such that three-parameters fits could reproduce the results of two-parameters fits presented in the literature. This experiment is therefore a crucial step in the way to the measurement of the two PMNS unknowns, if the precision goals can indeed be met.
The same kind of analysis we presented here must be clearly repeated at all of the proposed setups, something that goes beyond the scope of this paper but is extremely important to establish on solid grounds the quest for $`\theta _{13}`$ and leptonic CP violation in the near future.
## Acknowledgements
We would like to thank E. Fernรกndez-Martรญnez for extremely useful discussions and comments, D. Autiero, Y. Declais, B. Gavela, J.J. Gomez-Cadenas, P. Hernandez, P. Huber, P. Lipari, E. Lisi, M. Lusignoli, O. Mena, P. Migliozzi, T. Schwetz and W. Winter for discussions. The authors acknowledge the financial support of MEC through project FPA2003-04597, of CICYT-INFN through the โNeutrinos and others windows to new physics beyond the SMโ agreement and of the European Union through the networking activity BENE and the RTN European Program MRTN-CT-2004-503369.
## Appendix A
In this Appendix we review some of the statistical methods that have been proposed in the literature to take into account uncertainties on the external parameters in the measurement of ($`\theta _{13},\delta `$).
1) Comparison between naive two- and three-parameters fit
The obvious difference is in the CL contours that can be drawn in the two cases: for two-parameters fit the 90% CL corresponds to $`\mathrm{\Delta }\chi ^2=4.61`$, whereas for three-parameters fit is $`\mathrm{\Delta }\chi ^2=6.25`$. As a consequence, when a single minimum is found the projection of a three-parameters fit onto the two-dimensional contour is, in general, a bit larger. The second, not obvious, difference resides in that in the three-parameters fit the three-dimensional manifold automatically allows for a displacement of the clones solutions arranging for a lower $`\chi ^2`$ at the relative minima. This is indeed the case for the clones corresponding to wrong choices of $`s_{atm}`$ and of the $`\theta _{23}`$-octant (see , also). If the clones location moves in the three-dimensional manifold, the resulting projection onto the plane can be much larger than the two-parameters contour. This is indeed the main result of this paper and is discussed at length in Sect. 6.
2) Inclusion of a covariance matrix in the two-parameters fit
A fixed error range for any non-fitted parameter can be taken into account introducing a covariance matrix in the $`\chi ^2`$ function as follows:
$`\chi _{\{\overline{\alpha }\}}^2(\theta _{13},\delta )`$ $`=`$ $`{\displaystyle \underset{i,j}{}}\left\{\left[N_i(\theta _{13},\delta )N_i(\overline{\theta }_{13},\overline{\delta })\right]C_{ij}^1\left[N_j(\theta _{13},\delta )N_j(\overline{\theta }_{13},\overline{\delta })\right]\right\}_{\{\overline{\alpha }\}}`$
where $`C_{ij}`$ is the covariance matrix, $`i,j`$ refer to different channels at the same experiment or to different experiments and $`\{\overline{\alpha }\}`$ is a given set of external parameters. If the errors on the entries $`i`$ and $`j`$ of the covariance matrix are statistically independent, $`C`$ is
$$C_{ij}=\delta _{ij}\delta N_i^2+\underset{\alpha =1}{\overset{N_\alpha }{}}\frac{N_i}{\alpha }\frac{N_j}{\alpha }\sigma ^2(\alpha )$$
(14)
where $`\sigma (\alpha )`$ is the 1$`\sigma `$ error on the parameter $`\alpha `$. This procedure, followed in for the Neutrino Factory and in for the facilities considered in this paper, reproduces the enlargement of the two-parameters CL contours observed from a multi-parameter fit projected onto the ($`\theta _{13},\delta `$) plane. However, within this approach, the clones locations are not free to move in the multi-dimensional manifold to arrange for a lower $`\chi ^2`$: they are indeed stucked to the location in the ($`\theta _{13},\delta `$) plane that can be computed once known the external (fixed) parameters (see , again, and , Sect. 3.3). The displacement of the relative minima is indeed the characteristic feature of the multi-parameter fit, where $`N_\alpha `$ external parameters cooperate with $`\theta _{13}`$ and $`\delta `$ to locate lower $`\chi ^2`$ regions than those found in a two-parameters fit with fixed external parameters. The resulting regions are therefore large than those obtained including the covariance matrix in a two-parameters fit.
3) CP-coverage and marginalization over $`N_\alpha `$ external parameters
A parameter that can be used to compare in a condensed way the capability of different setups to measure the CP-violating phase $`\delta `$ has been proposed in . The CP-coverage is defined as follows:
$$\xi (\overline{\delta })=\mathrm{Coverage}\mathrm{in}\delta =\frac{1}{2\pi }U_{I=1}^{N_{deg}}\mathrm{\Delta }_I(\overline{\delta }),$$
(15)
is the fraction of the $`\delta `$-parameter space (i.e., $`2\pi `$) that is allowed at a given CL as a result of a measure when the input parameter is $`\overline{\delta }`$. The sum goes over $`N_{deg}`$ possible allowed regions induced by parameter degeneracies, each of them spanning an interval $`\mathrm{\Delta }_I(\overline{\delta })`$ of the parameter space. We take the union of these intervals to take into account possible partial overlaps of the $`\mathrm{\Delta }_I(\overline{\delta })`$. The smaller the CP-coverage, the better the measurement of $`\delta `$ at a given experiment. In particular, to distinguish a maximally CP-violating signal (i.e., $`\overline{\delta }=\pm 90^{}`$) at a given experiment from $`\delta =0^{}`$ or $`\delta =180^{}`$, the CP-coverage must be less than 0.5.
The definition of the CP-coverage above is, however, incomplete: we must still define how the $`\mathrm{\Delta }_I(\overline{\delta })`$ intervals are computed and which is the dependence of $`\xi (\overline{\delta })`$ on other parameters. Indeed, if all the parameters of the PMNS matrix but $`\delta `$ were measured, $`\xi (\overline{\delta })`$ would be just an involute way to express the expected error of an experiment for a certain value of $`\overline{\delta }`$. This has been called the โCP-patternโ, see Fig. 4(right) in Ref. . If, on the other hand, $`\theta _{13}`$ is also an unknown parameter, we can think of defining a function $`\xi (\overline{\theta }_{13},\overline{\delta })`$ and to plot it as a function of different values of $`\overline{\theta }_{13}`$ (called โCP-scalingโ in Ref. ) for a fixed value of $`\overline{\delta }`$. In this second case, for the particular value $`\overline{\delta }=0^{}`$, the โCP-scalingโ as a function of $`\overline{\theta }_{13}`$ is nothing else that the CP-sensitivity. For example, Fig. 5 in Ref. or Fig. 11 of Ref. represent the sensitivity to $`\delta `$ for varying $`\overline{\theta }_{13}`$ defined as a one-parameter fit where all mixing parameters have been measured but $`\delta `$ and $`\theta _{13}`$. The same idea is presented in Fig. 6 of Ref. and Fig. 13 of Ref. , where the plot represents the capability to distinguish a non-vanishing $`\delta `$ from $`\delta =0^{}`$ or $`\delta =180^{}`$ at a given one-parameter CL<sup>7</sup><sup>7</sup>7It should be stressed that it is not completely correct to present this โCP discovery potentialโ with one-parameter CL contours: being $`\theta _{13}`$ a parameter to be fitted together with $`\delta `$, we should consider two-parameters CL contours, instead. In this case the CP discovery potential can be smaller than when only $`\delta `$ is left as a free parameter, as a result of the fact that the two-parameters $`\chi ^2(\theta _{13},\delta )`$ can be lower than $`\chi ^2(\overline{\theta }_{13},\delta )`$ for specific values of $`\theta _{13}\overline{\theta }_{13}`$. This is indeed what reported in Fig. 7 of Ref. ..
To take into account the fact that, in general, the parameters of the mixing matrix are known only with a finite precision and that $`\theta _{13}`$ is completely unknown at present (and thus a one-parameter $`\delta `$-sensitivity plot is generally overestimating the performance of a given experiment), the authors of Ref. have proposed to compute the CP-coverage $`\xi (\overline{\delta })`$ by first minimizing a ($`N_\alpha +2`$)-parameter $`\chi ^2`$ over $`N_\alpha `$ external parameters. In this way, for any input pair $`\overline{\theta }_{13}`$ and $`\overline{\delta }`$, a two-dimensional surface of the $`\chi ^2`$ minimum as a function of $`\theta _{13}`$ and $`\delta `$ is generated. If we then minimize the resulting function in $`\theta _{13}`$, also, we can deduce a one-dimensional function of $`\delta `$ and of the input parameters representing the minimum value of the ($`N_\alpha +2`$)-dimensional $`\chi ^2`$ for a given value of $`\overline{\delta }`$. From this marginalized $`\chi ^2`$ we can finally compute the allowed $`\mathrm{\Delta }_I(\overline{\delta })`$ intervals imposing that $`\chi _{min}^2(\delta ,\overline{\delta })`$ be equal to a given one-parameter CL.
Clearly, this procedure can fail when multiple minima are present at each marginalization step. When multiple minima are present, the minimization procedure will generally look for the absolute minimum. In this way, the information on other relative minima in the $`\chi ^2`$ can be lost. This is particularly problematic since we know that, due to parametric degeneracies, multiple minima should be present. Marginalization near a second minimum will give a second function $`\chi _{min}^2(\delta ,\overline{\delta })`$, from which a new set of $`\mathrm{\Delta }_I(\overline{\delta })`$ intervals can be found. The procedure suggested in is indeed to marginalize around each of the expected minima in the multi-dimensional $`\chi ^2`$ and to draw several distinct one-dimensional functions $`\chi _{min}^2(\delta ,\overline{\delta })`$. Imposing on any of them the constraint $`\chi _{min}^2(\delta ,\overline{\delta })=CL`$, the full set of allowed regions in $`\delta `$ at a given one-parameter CL for a fixed input $`\overline{\delta }`$ is deduced and the $`\xi (\overline{\delta })`$ can be finally computed. Since the minimization procedure must be repeated several times (once per expected minima), it is useful to choose the starting point of the minimization algorithm in a clever way. In it is suggested to solve eqs. (2)-(5) as it has been done in Ref. to find the expected clone locations and to use the latter as starting points for the minimization. Applying the previous algorithm it is possible to deduce a $`\xi (\overline{\delta },\overline{\theta }_{13})`$ parameter that reduce significantly the overestimation of the experiment performances in the measurement of $`\delta `$ that is typical of the one-parameter $`\delta `$-sensitivity plots. We should therefore compare this procedure with the projection onto the $`\delta `$-axis of our three-parameters fits to understand if a residual overestimation is still present. To this purpose, in Fig. 11 we present the CP-coverage $`\xi (\overline{\theta }_{13},\overline{\delta })`$ and the projection onto the $`\delta `$-axis of the three-parameters CL contours for different values of $`\overline{\theta }_{13}`$ and $`\overline{\delta }`$. The results in the figure have been obtained using the considered $`\beta `$-Beam facility, for definiteness.
As it can be seen in the figure, some underestimation of the error on $`\delta `$ at the considered experiment is still present when computing $`\xi (\overline{\delta })`$ and comparing it with the $`\delta `$-axis projection of the three-parameters fits. The main interest of the CP-coverage parameter is that the algorithm described above can be iterated for any value of $`\overline{\theta }_{13}`$ to obtain a โCP-scalingโ for any given value of $`\overline{\delta }`$, thus replacing the (poorly reliable) one-parameter $`\delta `$-sensitivity plots. This will give the general picture of the expected error on $`\overline{\delta }`$ at a given facility, to be complemented in our opinion with multi-parameter fits to particular values of $`\theta _{13}`$ and $`\delta `$ to get a robust estimate of the facility performance.
## Appendix B: three-parameters fits
In this Appendix we present the projection of the three-dimensional 90% CL contours onto the ($`\theta _{13},\delta `$) plane for the three reference setups (the low-$`\gamma `$ $`\beta `$-Beam, the SPL Super-Beam and the 50 GeV Neutrino Factory) for different choices of the input pair ($`\overline{\theta }_{13},\overline{\delta }`$). The external input parameters in all plots are: $`\theta _{12}=32^{}`$, $`\mathrm{\Delta }m_{12}^2=8.2\times 10^5`$ eV<sup>2</sup>; $`\theta _{23}=40^{},\mathrm{\Delta }m_{23}^2=2.5\times 10^3`$ eV<sup>2</sup>. In this case, all choices of the two discrete variables $`s_{atm}`$ and $`s_{oct}`$ are presented together and no comparison with two-parameters contours is shown.
For the $`\beta `$-Beam plots we compare the impact of the present uncertainties on the atmospheric parameters (third column of Tab. 1) with that of the expected uncertainties after T2K-I (last column of Tab. 1): $`\theta _{23}[38^{},43^{}][48^{},52^{}]`$ and $`\mathrm{\Delta }m_{23}^2[2.42,2.61]\times 10^3`$ eV<sup>2</sup> for $`s_{atm}=+`$ and $`\mathrm{\Delta }m_{23}^2[2.46,2.64]\times 10^3`$ eV<sup>2</sup> for $`s_{atm}=`$, . Notice that the error on $`\theta _{23}`$ is just the expected error at T2K-I, , shifted around $`\theta _{23}=40^{}`$. It has been shown in Sect. 5.2 that the $`\nu _\mu `$ disappearance channel at the Super-Beam is rather effective in reducing the uncertainties on the atmospheric parameters (whereas the $`\nu _e`$ disappearance channel at the low-$`\gamma `$ $`\beta `$-Beam is useless to this purpose, see Ref. ). In this case we therefore do not present results using โpresentโ and โexpectedโ uncertainties, but we just combine the results from the appearance and disappearance channel. Finally, for the Neutrino Factory we have considered only the two appearance channels $`\nu _e\nu _\mu ,\nu _\tau `$ with the atmospheric parameters with the expected uncertainties after T2K-I.
In general, a โpessimisticโ systematic error, $`ฯต^\pm =5`$%, has been used in appearance channels. A 2% systematic error has been used in disappearance channels. However, we checked that using a โpessimisticโ 5% systematic error in the disappearance channel does not change significantly our results.
The input parameters are: $`\overline{\theta }_{13}=2^{},7^{}`$; $`\overline{\delta }=90^{},45^{},0^{},45^{},90^{}`$.
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# Spin and temperature dependent study of exchange and correlation in thick two-dimensional electron layers.
## I Introduction.
The 2D electron systems (2DES) present in GaAs or Si/SiO<sub>2</sub> structures access a wide range of electron densities, providing a wealth of experimental observationskrav . The 2D electrons reside in the x-y plane and also have a transverse ($`z`$-dependent) density $`n(z)`$ which is confined to the lowest subband of the hetero-structureafs . The 2D character arises since the higher subbands are sufficiently above the Fermi energy $`E_F`$. Then the physics depends only on the โcoupling parameterโ $`\mathrm{\Gamma }`$ = (potential energy)/(kinetic energy), the $`z`$-distribution $`n(z)`$, the spin-polarization $`\zeta `$, and the temperature $`T`$ which has to be significantly smaller than the Fermi energy $`E_F`$ to preserve the 2-D character. The $`\mathrm{\Gamma }`$ for the 2DES at the density $`n`$ is equal to the mean-disk radius $`r_s=(\pi n)^{1/2}`$ per electron, expressed in effective atomic units which depend on the bandstructure mass $`m_b`$ and the โbackgroundโ dielectric constant $`ฯต_b`$. The $`z`$-motion in the lowest subband may have widths of $`600`$ ร
, in GaAs when $`r_s`$ is $`6`$, in heterojunction insulated-gate field-effect transistors (HIGFET) which have been an object of recent studieszhu . Similarly, other nanostructures (e.g., quantum dots) contain electrons confined to a micro-region in the $`x`$-$`y`$ plane, and have a sizable $`z`$ extensionwilliams . Hence layer-thickness effects are important in many areas of nanostructure physics. Appropriate correlation functionalsmartin for such systems are still unavailable, even though the exchange functional for Fang-Howard distributions is known sternjjap .
Layer thickness effects are a long standing probe of exchange and correlation theories in 3D electron slabsperdew . The relevance of the finite size of the 2D layers had also been considered within the quantum Hall effectahm ; dassarmaqh , and more recently in the context of the $`g`$-factor and the effective mass $`m^{}`$ of 2-D layertutuc ; tan ; morsen ; asgari ; zhang ; cdw . In the early days of the application of diagrammatic perturbation theory (PT) to condensed-matter problems, it was normal to attempt to calculate various many-body properties like the effective mass $`m^{}`$, the effective $`g`$-factor $`g^{}`$, and corrections to the total energy using perturbation methods. The need to go beyond the random-phase approximation (RPA) was rapidly appreciated and lead to the work of Hubbard, Rice, Vosko, Geldart and othersgeldart . The common experience with the generalized RPA method, when applied to 3D electrons and to ideally thin 2D layer is that it predicts a spin-polarized phase at unrealistic high densities (low $`r_s`$), while Quantum Monte Carlo (QMC) simulations suggest a density near $`r_s`$ 25-26 in 2D. RPA methods lead to negative pair-distribution functions (PDFs), incorrect local-field corrections in the response functions, and disagreement with the compressibility sum rule and other formal conditions. Recent attempts to calculate the effective mass $`m^{}`$ for ideally thin layersasgari ; zhang show strong disagreement with the $`m^{}`$ obtained from QMC simulationskwon . In fact, the main thrust of the programs of Singwi, Tosi et al. (STLS)stls , and Ichimaru et al.ichimaru was to overcome the shortcomings of the RPA-like approach via non-perturbative methods. Many-body calculations where different parts of the calculation (e.g, vertex corrections, local-field corrections, etc.) are obtained from different sources, (e.g, vertex corrections from a model, local-field corrections or correlation energies from QMC, and some other quantities fitted to sum rules etc.) and combined together, have also appeared. Unless the same โmixtureโ is used to calculate a multitude of properties and shown to lead to consistent results, such theories have to be treated with caution.
We have introduced an approximate method for strongly correlated quantum systems where the objective is to work with the PDF of the quantum fluid, generated from an equivalent classical Coulomb fluid whose temperature $`T_q`$ is chosen to reproduce the correlation energy of the original quantum fluid at $`T=0`$. The classical PDFs are calculated via the hyper-netted-chain (HNC) equation, and the method is called the CHNC. As the method has been described in previous publicationsprl1 ; prb ; prl2 ; prl3 , and successfully applied to a variety of problemshyd ; 2valley ; locf ; eplmass , only a brief account is given here, mainly to help the reader. The PDFs are obtained from an integral equation which can be recast into a classical Kohn-Sham form where the correlation effects are captured as a sum of HNC diagrams and bridge diagrams.
$$g_{ij}(r)=\mathrm{exp}^{\beta \left(P_{ij}(r)+V_{cou}(r)+N_{ij}(r)+B_{ij}(r)\right)}$$
(1)
The temperature of the classical fluid $`T_q=1/\beta `$ is chosen such that at $`T=0`$ the calculated classical $`g(r)`$ recovers the known correlation energy of the fully spin-polarized 2D electron fluid at the given density. This fitting has been done in ref. prl2 , and $`T_q`$ is known as a function of $`r_s`$. At finite-$`T`$, the classical fluid temperature $`T_{cf}`$ is taken to be
$$T_{cf}=(T_q^2+T^2)^{1/2}$$
(2)
This $`T_{cf}`$ is used for all spin polarizations. In Eq. 1 the indices $`i,j`$ label the spin states. The $`P_{ij}(r)`$ is chosen to ensure that $`g_{ij}(r)`$ reduces to the explicitly known non-interacting PDF, $`g_{ij}^0(r)`$, when the Coulomb interaction $`V_{cou}(r)`$ is switched off. Thus $`P_{ij}(r)`$ takes care of Pauli-exclusion effects and ensures that the โFermi-holeโ is exactly recoveredlado . Also, $`N_{ij}(r)`$ is a sum of terms that appear in the hyper-netted-chain equation, while $`B_{ij}(r)`$ contains 3-body and higher-order diagrams not contained in the nodal term $`N_{ij}(r)`$. The latter depends implicitly on $`g_{ij}(r)`$, and is evaluated via the Ornstein-Zernike equation. The bridge term is very difficult to evaluate, but the hard-sphere fluid provides a good approximation. That is, in 2D, we specify $`B_{ij}(r)`$ by specifying the hard-disk radius $`r_H`$, or equivalently the packing fraction $`\eta `$, and use the Percus-Yevik approachharddisk . As the parallel-spin three-body clusters are suppressed by the Pauli exclusion, we use only a single bridge function, viz., $`B_{12}(r)`$. This makes the bridge interaction effectively independent of $`\zeta `$. Bulutay and Tanatar, and also Khanh and Totsujibuluty , have examined variants of CHNC without a bridge function (this is some what equivalent to neglecting back-flow terms in QMC simulations of 2D systems). The hard-sphere radius $`r_H`$, and the packing fraction $`\eta `$ are given byprl2 :
$`r_H`$ $`=`$ $`2r_s\eta ^{1/2}`$ (3)
$`\eta `$ $`=`$ $`0.1175r_s^{1/3}(t_{cf}+t^2/2)`$ (4)
$`t_{cf}`$ $`=`$ $`T_{cf}/E_F,t=T/E_F`$ (5)
A plot of the bridge function for a few typical cases is given in Fig. 1.
The CHNC method was applied to the 3D and 2D electron fluidsprl1 ; prb ; prl2 ; prl3 , to dense hydrogen fluidhyd , and also to the two-valley system in Si/SiO<sub>2</sub> 2DES2valley ; eplmass . In each case we showed that the PDFs, energies and other properties obtained from CHNC were in excellent agreement with comparable QMC results. The advantage of the CHNC method is that it affords a simple, semi-analytic theory for strongly correlated systems where QMC becomes prohibitive or technically impossible to carry out. The classical-fluid model allows for physically motivated treatments of complex issues like three-body clustering etc., which are difficult within quantum methods. The disadvantage, typical of such many-body approaches, is that at present it remains an โextrapolationโ taking off from the results of a model fluid. The fully spin-polarized infinitely-thin uniform 2DES is the model fluidprl2 , as the $`E_c`$ of this one-component system is accurately known.
In this study we use a method of replacing the inhomogeneous electron distribution $`n(\stackrel{}{r},z)=n(\stackrel{}{r})n(z)`$ by a homogeneous density systemggsavin ; cdw i.e., a constant-density approximation (CDA) where the transverse density $`n(z)`$ is a constant within a slab of width $`w`$, and zero outside. The CDA avoids difficulties associated with gradient expansions noted in ref.martin . Also, the CDA presents a unified approach to quasi-2D distributions like the Fang-Howard modelafs , the quantum-well model etc. $`E_{xc}`$ for such distributions is calculated as a function of the spin polarization $`\zeta `$, 2-D density parameter $`r_s`$, the CDA width $`w`$ and the temperature $`T`$. This enables us to determine physical quantities related to the Landau Fermi liquid parameters. Thus the spin susceptibility enhancement, the effective mass $`m^{}`$, and the Landรฉ $`g`$ factor for the quasi 2D electrons are presented.
## II The quasi-2D interaction
The transverse distribution $`n(z)`$ of the 2D electrons is given by the square of the lowest subband wavefunction $`\varphi (z)`$ of the heterostructure, calculated within the envelope approximation. The nature of the materials used (e.g, Si/SiO<sub>2</sub> or GaAs) and the doping profiles determine the confining potential and the electron density $`n(z)`$ in the z-direction. Typically, $`n(z)`$ may be modeled by one of the following forms.
$`n(z)`$ $`=`$ $`\delta (z),\text{ideal thin layer}`$ (6)
$`=`$ $`(2/w)\mathrm{sin}^2(\pi z/w),\text{infinite well}`$ (7)
$`=`$ $`(b^3/2)z^2\mathrm{exp}(bz),\text{Fang-Howard}`$ (8)
$`=`$ $`1/w,|z|w\text{constant-density model}`$ (9)
The second and third are frequently used approximate (but adequate) models, while we present the fourth model, the CDM. This is an excellent constant-density approximation (to be called the constant-density approximation, CDA) to generate the effective 2-D potential arising from most $`n(z)`$ distributions, on suitably choosing $`w`$. The Fang-Howard (FH) formafs ; fangh , contains the parameter $`b`$, and is normalized in the range $`0`$ to $`\mathrm{}`$. The FH parameter $`b`$ is such that $`b^3=(48\pi /a_0)(n_d+11n_s/32)`$, where $`n_s`$ is the 2D electron density $`n`$, and $`n_d`$ is the depletion charge densityafs . The material parameters are contained in the effective Bohr radius $`a_0=ฯต_b\mathrm{}^2/m_be^2`$ defined in terms of the usual constants, $`ฯต_b`$ and $`m_b`$ being the dielectric constant and the band mass. For the devices used in ref. tan ; zhu , the depletion density has been reported to be negligiblemorsen . Then $`b^3=33/(2r_s^2)`$, in atomic units.
### II.1 The constant-density model.
We denote the Coulomb potential in an infinitely thin layer by $`V(r)=1/r`$, while $`W(r)`$ is used for the effective 2-D potential of a thick layer. The effective 2D-Coulomb potential $`W(r)`$ between two electrons having coordinates ($`\stackrel{}{r_1},z_1`$) and ($`\stackrel{}{r_2},z_2`$), with $`\stackrel{}{r}=\stackrel{}{r_1}\stackrel{}{r_2}`$ is given by,
$$W(r)=_0^{z_m}_0^{z_m}\frac{dz_1dz_2n(z_1)n(z_2)}{[r^2+(z_1z_2)^2]^{1/2}}$$
(10)
Here $`z_m`$ is $`\mathrm{}`$ for FH, while $`z_m=w`$ for the others. The potential $`W(r)=(1/r)F(r)`$ and the form factor $`F(r)`$ reflects the effect of the $`z`$-extension of the density. There is no analytic form for $`F(r)`$ in the Fang-Howard case, while the $`q`$-space form, $`F(q)`$ is knownafs . If the dielectric constants of the barrier and well material were assumed equal, then the Fang-Howard form factor is:
$$F(q)=[1+\frac{9q}{8b}+\frac{3q^2}{8b^2}][1+\frac{q}{b}]^2$$
(11)
Here we derive a potential $`W(r,w)`$ for the constant-density model(CDM) which is, to an excellent approximation electrostatically equivalent to the the 2D potential for any reasonable $`n(z)`$. These FH-type distributions are themselves convenient fits to the self-consistent Schrodinger solutions and have uncertainties of a few percent. The CDA holds well within such limits. The potentials are explicitly shown in Fig. 2 for the FH form where we have taken an extreme example with $`b=0.1`$. The method of replacing an inhomogeneous distribution by a uniform distribution is suggested by the observation that the non-interacting total pair-correlation function $`h^0(r)`$ has the form $`n(r)^2`$, where $`n(r)`$ is the density -profile around the Fermi hole. In our case we wish to replace the inhomogeneous $`n(z)`$ by a constant-density distribution $`n_{cd}`$ which has the same electrostatic potential in the 2-D plane as $`n(z)`$.
$$n_{cd}=1/w=n(z)^2๐z$$
(12)
Since the subband distribution is normalized to unity, the width $`w`$ of the constant-density slab is simply $`1/n_{cd}`$. Starting from different objectives, Gori-Giorgi et al.ggsavin have already proposed this formula for determining an average density for an inhomogeneous density, in the context of pair-distributions in 2-electron atoms. We have also shown the utility of the CDA in estimating the correlation energy of the 2DES in the Wigner-crystal phasejost . The $`w`$ of the CDA is somewhat different from the โthicknessโ $`3/b`$ often assigned to the FH distribution. In fact, the constant-density slab width $`w`$ for the Fang-Howard $`b`$ is given by $`w=16/(3b)`$.
The quasi-2D potential for a CDM of width $`w`$ is given by
$`W(r)`$ $`=`$ $`V(r)F(s),s=r/w,t=\sqrt{(1+s^2)}`$ (13)
$`F(s)`$ $`=`$ $`2s\left[\mathrm{log}{\displaystyle \frac{1t}{s}}+1t\right]`$ (14)
This potential tends to $`1/r`$ for large $`r`$, and behaves as
$$\frac{2}{w}\left(\mathrm{ln}\frac{2w}{r}+\frac{r}{w}1\right)$$
for $`r<w`$. Thus the short-range behaviour is logarithmic and weaker than the Coulomb potential. The $`k`$-space form of the CDM potential is:
$`V_{usm}(k,w)`$ $`=`$ $`V(k)F(p),p=kw`$ (15)
$`F(p)`$ $`=`$ $`(2/p)\{(e^p1)/p+1\}`$ (16)
The form factors $`F(s)`$ and $`F(p)`$ tend to unity as $`w0`$. These $`r`$-space and $`k`$-space analytic forms of the CDM lead to analytic formulae for the FH form. In our work we assume that a given distribution has been replaced by an equivalent uniform-slab distribution, and only the final $`W(r)`$ potential enters into the exchange-correlation calculations (the numerical work has been checked via direct calculations as well). In the case of GaAs-HIGFETS, if $`n_d`$ could be neglected, the $`r_s`$ parameter specifies the $`b`$ parameter and hence the width $`w`$ of the CDM. Then $`br_s^{2/3}`$ and $`w=2.09494r_s^{2/3}`$.
## III Exchange Free Energy for quasi-2D layers.
The exact exchange free energy $`F_x`$ for 2D layers of finite thickness can be readily evaluated using the quasi-2D potential $`W(r)`$ and the noninteracting pair-distribution functions $`g_{ij}^0(r)`$ of the 2-D fluid. Only the parallel-spin case $`i=j`$ is relevant. Also, $`g_{ij}^0(r)`$ for a slab of finite thickness is identical to that for an ideally thin 2D layer, both at $`T=0`$ and at finite-$`T`$. In fact, we find that the $`T`$ dependence of the $`F_x`$ of layers of finite thickness is very close to that of the ideally thin case.
### III.1 Ideally-thin layer
The first-order unscreened exchange free energy $`F_x`$ consists of $`F_x^i`$, where $`i`$ denotes the two spin species. At $`T=0`$ these reduce to the exchange energies:
$$E_i^x/n=\frac{8}{3\sqrt{\pi }}n_i^{1/2}$$
(17)
Here $`n_1=n(1+\zeta )/2`$, and $`n_2=n(1\zeta )/2`$. Then the exchange energy per particle at $`T=0`$, i.e., the total $`E_x/n`$ becomes
$$E_x/n=(E_1^x+E_2^x)/n=\frac{8}{3\pi r_s}[c_1^{3/2}+c_2^{3/2}]$$
(18)
where $`c_1`$ and $`c_2`$ are the fractional compositions $`(1\pm \zeta )/2`$ of the two spin species.
We define a reduced temperature $`t=T/E_F`$, $`E_F=\pi n`$, and the species-dependent reduced chemical potentials $`\mu _i^0/T`$ by $`\eta _i`$, reduced temperatures $`t_1=t/(1+\zeta )`$ and $`t_2=t/(1\zeta )`$, based on the two Fermi energies $`E_{F1}`$ and $`E_{F2}`$ which are $`E_F(1\pm \zeta )`$. Then we have:
$$F_i^x/E_i^x=\frac{3}{16}t_i^{3/2}_{\mathrm{}}^{\eta _i}\frac{I_{1/2}^2(u)du}{(\eta _iu)^{1/2}}$$
(19)
The $`I_{1/2}`$ is the Fermi integral defined as usual:
$$I_\nu (z)=_0^{\mathrm{}}\frac{dxx^\nu }{1+e^{xz}}$$
(20)
The $`\eta _i`$ are given by
$$\eta _i=\mathrm{log}(e^{1/t_i}1)$$
(21)
In the paramagnetic case Eq. 19 reduces to the result given by Isihara et al. (see their Eq.3.4; they use a different definition of the Fermi integral). For small values of $`t`$, the exchange energy is of the form,
$$E_x(r_s,t)=E_x(r_s,0)[1+(\pi ^2/16)t^2\mathrm{log}(t)0.56736t^2+\mathrm{}]$$
(22)
The total exchange free energy is $`F_x=\mathrm{\Sigma }F_i^x`$. The accurate numerical evaluation of Eq. 19 requires the removal of the square-root singularity by adding and subtracting, e.g, $`I^2(|\eta |)/(v|\eta |)^{1/2}`$ for the case where $`\eta `$ is negative, and $`v=u`$, and so on.
A real-space formulation of $`F_x`$ = $`F_1^x+F_2^x`$ using the zeroth-order PDFs fits naturally with the approach of our study. Thus
$$F_x/n=n\frac{2\pi rdr}{r}\underset{i<j}{}h_{ij}^0(r)$$
(23)
Here $`h_{ij}^0(r)=g_{ij}^0(r)1`$. In the non-interacting system at temperature $`T`$, the antiparallel $`h_{12}^0`$, viz., $`g_{12}^0(r,T)1`$, is zero while
$$h_{11}^0(๐ซ)=\frac{1}{n_i^2}\mathrm{\Sigma }_{๐ค_1,๐ค_2}n(k_1)n(k_2)e^{i(๐ค_1๐ค_2)๐ซ}=[f(r)]^2$$
Here k, r are 2-D vectors and $`n(k)`$ is the Fermi occupation number at the temperature $`T`$. At $`T=0`$ $`f(r)=2J_1(k_ir)/k_r`$ where $`J_1(x)`$ is a Bessel function. As a numerical check, we have evaluated the exchange free energy by both methods, i.e., via k-space and r-space formulations.
We present a convenient analytic fit to the exchange free energy which is a universal function $`F_x(t)/E_x`$, for arbitrary $`\zeta `$. That is, the same function applies to any component, on using the reduced Fermi temperature of the spin species. The total $`F_x`$ is obtained by adding both spin contributions. The analytic fit is:
$`F_i^x(t,\zeta )/E_i^x(\zeta )=`$
$`{\displaystyle \frac{1+C_1t_i+C_2t_i^2}{1+C_3t_i+C_4t_i^2}}\mathrm{tanh}(1/\sqrt{t_i})`$
The fit coefficients $`C_i`$ are 3.27603, 4.81484, 3.33100, 6.51436. The temperature $`t_i`$ is $`t/(1\pm \zeta )`$, appropriate to the spin polarization. The exchange effects in the 2DES decay more slowly with temperature than in the 3D case where a tanh$`(1/t)`$ factor appears in Eq. 3.2 of ref. pdw84, . The above form does not explicitly contain the low-temperature logarithmic term, but it reproduces the value of 0.99382 at $`t=0.05`$, while the numerical integration gives 0.9939497. Similarly, at $`t=1`$, 10 and 30 the fit (integral) returns 0.63839 (0.63839), 0.22999 (0.22990), and 0.13421 (0.13410) respectively.
### III.2 Thick 2-D layers.
The $`T=0`$ exchange energy is modified by the layer thickness $`w`$. The expression for $`g_{ij}^0(r)`$ depends only on $`x=r/r_s`$. Similarly, the quasi-potential $`W(r)`$ depends on the reduced variable $`s=r/w`$. Hence the exchange energy of the quasi-2D layer depends only on the ratio $`w_s=w/r_s`$. The exchange energy per electron at a density $`n`$, given by $`r_s`$, polarization $`\zeta `$, in a layer of width $`w`$ is given by:
$$Ei_x(r_s,\zeta ,w)=\frac{1}{2}nr_s_0^{\mathrm{}}2\pi ๐xh^0(x,\zeta )F(x/w_s)$$
(25)
Even though we have analytic forms for $`W(r)`$ and $`h^0(r)`$ as well as their Fourier transforms, we have not found a convenient analytic result for the exchange energy at $`T=0`$. However, the results can be parametrized by simple analytic-fit formulae:
$`E_x(r_s,\zeta ,w)`$ $`=`$ $`E_x(r_s,\zeta ,0)Q(w_s,\zeta ),w_s=w/r_s`$
$`Q(w_s,\zeta )`$ $`=`$ $`{\displaystyle \frac{1+A(\zeta )\sqrt{w_s}}{1+B(\zeta )\sqrt{w_s}+C(\zeta )w_s}}`$
Here $`E_x(r_s,\zeta ,0)`$ is the exchange energy of the ideally thin system given by eq. 18. The ratio $`Q(w_s,\zeta )`$ is a measure of the reduction in the exchange energy due to the thickness effect. Since the effect depends on $`w_s=w/r_s`$, for a given thickness, the effect is greater for high density samples. If the depletion density $`n_d`$ in HIGFETS could be neglected, and if the exchange-correlation energy $`E_{xc}`$ is not included in the energy minimization which determines the Fang-Howard parameter $`b`$, then $`w_s2.09r_s^{1/3}`$. The inclusion of $`E_{xc}`$ in self-consistently determining $`b`$ changes $`b`$ by $`2\%`$ for low $`r_s`$, but the effect becomes less important at higher $`r_s`$. At $`r_s`$=1, and 10 for $`\zeta =0`$, the ratio $`Q`$ is 0.652 and 0.794 respectively. The reduction from the ideally-thin 2D form is clearly substantial. The exchange free energy $`F_x(r_s,\zeta ,T,w)`$ at finite-$`T`$, for layers with thickness $`w`$ is found to be adequately approximated by the temperature factor of the ideally-thin system. However, in calculating the effective mass from the finite-$`T`$ energies, we make independent calculations near $`T=0`$ and dอกo not use the fit given here.
The exchange energy for a HIGFET with $`n_d=0`$ is shown in Fig. 3 as a function of $`r_s`$ for $`\zeta =0`$ and temperature $`T/E_F`$=0 and 0.2.
## IV The exchange-correlation energy for quasi-2D layers.
The correlation function $`h^0(r)`$ yields exact exchange energies for arbitrary layer thicknesses. In contrast, the correlation energies require a coupling-constant integration over the pair functions $`g(r,\zeta ,w,\lambda )`$ calculated using the quasi-2D potential $`\lambda W(r)`$ for each $`\lambda `$. These $`g(r,\zeta ,w)`$ can be calculated using the CHNC. On the other hand, the unperturbed-$`g`$ approximation, found to be useful in Quantum Hall effect studiesahm , has been exploited by De Palo et al.morsen . They have used the pair functions $`g(r,\zeta ,w=0)`$ of the ideally thin layer obtained form QMC to calculate a correction energy $`\mathrm{\Delta }`$ given by,
$$\mathrm{\Delta }E=(n/2)2\pi r๐r[W(r)V(r)]h(r,\zeta ,w=0)$$
(26)
Then the total exchange-correlation energy $`E_{xc}(r_s,\zeta ,w)`$ is obtained by adding to $`\mathrm{\Delta }`$ the known $`E_{xc}`$ of the ideally-thin system. The above equation can also be applied at finite temperatures using the finite-$`T`$ pair functions $`g(r_s,\zeta ,T)`$ obtainable from the CHNC procedure.
De Palo et almorsen have performed Diffusion-Monte-Carlo simulations at $`r_s`$ =5 for HIGFETS with $`b=0.8707`$, i.e, a CDM width $`w`$=6.1256 a.u., and find that the error in this approach compared to the full simulation is about 2%. This full QMC result at $`r_s=5`$ is shown in the lower panel of Fig. 3. Since the HIGFET system approximates to a thin-layer as $`r_s`$ increases, this approach is probably satisfactory for $`r_s`$ 5. The method becomes unreliable for small $`r_s`$, and definitely fails below $`r_s=2`$. Also, the โunperturbed-$`g`$โ approximation fails to include the renormalization of the kinetic energy picked up via the coupling constant integration over the fully consistent $`g(r,w)`$. We report results(Fig. 3) from the full coupling-constant integration of the $`g(r,w)`$, (Fig. 3, lower panel, CHNC) as well as from the โunperturbed-$`g`$โ approximationmorsen used by De Palo et al.
In parametrizing the quasi-2D correlation energy $`E_c(r_s,w)`$, we present an intuitive model of $`E_c(r_s,w)`$. For small $`r_s`$, the ratio $`w/r_s`$ is large and the electrons are like 1-D wires with the axis normal to the 2D plane and interacting with a $`\mathrm{log}(w/r)`$ interaction (c.f., Eq. 13). However, at large $`r_s`$ we have 3-D like electron disks with $`w`$ and $`r_s`$ of comparable magnitude in the density regimes of interest in HIGFETS. Thus we model the quasi-2D $`E_c`$ as an interpolation between a 1D like form and a 3D like form. First we consider a purely 3D model. Given the 2D-density $`r_s`$ and an effective CDM width $`w`$, we define an effective 3D density parameter $`r_s^{3d}`$, purely for calculating its correlation energy. It will be seen that this 3D model is excellent for $`r_s>7`$. When $`r_s`$ becomes small (i.e, less than $`3`$), the effective width of the 2D layer, viz., $`w/r_s`$ becomes large and a 1D log-interaction modelsamaj is needed. To capture the rod-like regime, we define the โrod likeโ correction $`\mathrm{\Delta }E_c`$ for $`r_s<7`$ by:
$$\mathrm{\Delta }E_c(\zeta =0)=a_0+a_L\mathrm{log}(1/r_s)+a_1r_s$$
(27)
where, for HIGFETS, $`a_L=0.0221788`$, $`a_1=0.00365169`$, and $`a_0=0.0192979`$ for $`\zeta =0`$. This $`\mathrm{\Delta }E_c`$ is added to the 3-D slab form given below. The fit parameters for the $`\zeta =1`$ โrod-likeโ correction are $`a_0`$=0.013337, $`a_L`$=0.0084787, and $`a_1`$=0.0006821, to be applied for $`r_s<15`$.
For the 3D slab-like regime (i.e, $`r_s>7`$ for $`\zeta `$=0, $`r_s>15`$ for $`\zeta =1`$) we define a $`\zeta `$-dependent 3-D density parameter and a correlation energy via:
$`E_c(r_s,\zeta ,w)`$ $`=`$ $`E_c^{3d}(r_s^{3d},\zeta )`$ (28)
$`R_s`$ $`=`$ $`r_s/F(r_s)`$ (29)
$`r_s^{3d}`$ $`=`$ $`[wR_s^2]^{1/3}Z(\zeta )`$ (30)
$`Z(\zeta )`$ $`=`$ $`{\displaystyle \frac{2\sqrt{2}}{(4\zeta )^{1/2}}}(c_1^{1.5}+c_2^{1.5})`$ (31)
The 3D correlation energy $`E_c^{3d}(r_s^{3d},\zeta )`$ is that given by, e.g., Ceperley and Aldercepalder , or Gori-Giorgi and Perdewg-gp . We see from the lower panel of Fig. 3 that the correlation energy for small $`r_s`$, calculated using the 3D slab begins to go below the โunperturbed-gโ approximation of Eq. 26, consistent with the trend of the ideal 2D gas, and the trend of the only QMC data point available for a HIGFET, at $`r_s=5`$. The exchange-correlation energy obtained from the full CHNC calculation is in excellent agreement with the QMC datum. The curve labeled โSlab+rodsโ in Fig. 3 is the combined formula, Eq. 27, and Eq. 28, for the correlation energy/electron, with, for example, the region $`r_s7`$ for $`\zeta =0`$ obtained by linear interpolation between $`r_s`$=6 and $`r_s=8`$. This is clearly seen to reproduce the full CHNC results very well.
### IV.1 Correlation energy at finite temperatures.
The correlation contribution to the Helmholtz free energy of the ideal 2-D layer, and layers of thickness $`w`$ can be easily calculated using the approach of Eq. 26, where the CHNC-generated finite-$`T`$ pair functions are used. A typical set of results at very low temperatures is given in Table 2. Here we have also given the packing fraction $`\eta `$ of the hard-sphere bridge function used to mimic the three-body and multi-body correlation contributions. As discussed in earlier workeplmass , the $`F_x`$ and the $`F_c`$ at very low $`T`$ contain logarithmic terms which cancel with each other, so that the sum $`F_{xc}=F_x+F_c`$ is free of such terms. From our numerical work we find that the $`T`$ dependence of the $`F_x`$ and $`F_c`$ of layers of finite thickness is very similar to that of ideally thin layers. Hence we assume that the logarithmic corrections are also similar. At $`r_s=5`$ the cancellation is good to about 75%, and this improves as $`r_s`$ increases. Although the two-component fluid (up spins and down spins) involves three distribution functions, we have, as beforeprl2 , used only one hard-disk bridge function, $`B_{12}`$, as clustering effects in $`g_{ii}`$ are suppressed by the Pauli-exclusion. However, at high densities (low $`r_s`$), the use of three bridge functions seems to be needed for satisfying various subtle features that are needed to ensure the exact cancellation of logarithmic energy terms etc. Instead of introducing additional features into the CHNC method, we have however retained the single bridge-function model that was used by us so farprl2 .
### IV.2 The transition to a spin-polarized phase.
QMC simulations as well as CHNC calculations show that there is a spin-polarization transition (SPT) in the ideally-thin 2D electron fluid near $`r_s26`$. On the other hand, the correlation contributions dominate over the exchange energy in the 2-valley 2D system in Si MOSFETS, and the SPT is suppressed2valley . The rapid increase in $`m^{}`$, with $`g^{}`$ remaining unchanged while $`r_s`$ is increased, observed by Shashkin et al.shash was found to be consistent with this pictureeplmass . In finite-thickness 2D layers, as the CDM width $`w`$ increases, the location of the SPT is pushed to higher values, as seen in Fig. 4. In the case of HIGFETS used by, e.g., Tan et al.tan , the width $`w`$ increases with $`r_s`$, but only as $`r_s^{2/3}`$. Thus at $`r_s`$=26, the $`w`$ is only 18.4, and the SPT remains intact and occurs at a somewhat higher $`r_s`$, as shown by de Palo et almorsen , and also in Fig. 4. A natural consequence of delaying the SPT is to decrease the spin-susceptibility enhancement. The effective thickness of the quasi-2D layer can be increased by suitably designing the shape of the potential well, or including an additional subband, and in this case the SPT can be circumvented. However, a discussion of higher subband effects is outside the scope of this study.
## V The spin-susceptibility, effective mass and the $`g`$-factor
The results for the exchange-correlation free energy $`F_{xc}(rs,\zeta ,T)`$ for the ideal 2D system and the thick-layer system contain all the information needed to calculate the spin-susceptibility enhancement, the effective mass $`m^{}`$ and the effective Landรฉ factor $`g^{}`$, for the ideal system and the thick layer. In fact, the following quantities are calculated from the respective second derivatives of the energy.
$`m^{}`$ $`=`$ $`C_v/C_v^0=1+{\displaystyle \frac{\left[^2F_{xc}(t)/t^2\right]}{\left[^2F_0(t)/t^2\right]}}`$ (32)
$`\chi _P/\chi _s`$ $`=`$ $`(m^{}g^{})^1=1+{\displaystyle \frac{\left[^2F_{xc}(\zeta )/\zeta ^2\right]}{\left[^2F_0(\zeta )/\zeta ^2\right]}}`$ (33)
We use the available QMC results for the ideal 2D exchange-correlation energy at $`T=0`$, and where convenient, the QMC pair-distribution functions at $`T=0`$ as parametrized by Giri-Giorgi et alggpair . The CHNC is used to obtain the pair-functions for situations (e.g, at finite-$`T`$ and finite thickness) where the QMC data are simply not available or difficult to use. In most cases, replacing the QMC-PDFs with the CHNC ones, or using the โunperturbed-$`g`$โ approximation leads to relatively small changes. The exception is in the calculation of $`m^{}`$, where the โunperturbed-$`g`$โ approximation, Eq. 26, is inadequate.
### V.1 The effective mass $`m^{}`$.
In Fig. 6 we present our results for the ideal-2DES. No experimental results are available for this case, but limited QMC simulationskwon as well as results from diagrammatic perturbation theory (PT) asgari ; zhang are available.
The CHNC based $`m^{}`$ values have an error of at most $`\pm `$ 2%. The ideal 2D-layer $`m^{}`$ shows a rapid rise around $`r_s`$=2 to 5, in good agreement with the four values from QMC, and then slows down in strong contrast to the $`m^{}`$ proposed by Asgari et al., and Zhang et al. We have displayed the model denoted $`G_+\&G_{}/D`$ by Asgari et al., as being their optimal choice from among the many models given in Ref. asgari , where they also contest the analysis of Zhang et al. The PT approaches are partly semi-phenomenological in that QMC data are input into local-field factors and other ingredients of the calculation; the choice of the vertex functions, treatment of on-shell or off-shell effects, whether to use self-consistent propagators or not, etc., are components of the prescription used by different workers. The strong disagreement of the PT-based $`m^{}`$ with the QMC-based $`m^{}`$ is notable.
The CHNC method has some similar uncertainties, especially in the use of a Percus-Yevik hard-sphere bridge function $`B(r)`$ to capture the back-flow like three-body contributions to the PDFs and the total energy. As seen in Fig. 1, the $`T`$ dependence of $`B(r)`$ seems quite small, and our initial calculations of $`m^{}`$, reported in Ref. ssc were based on the zero-$`T`$ form of $`B(r)`$. This leads to an $`m^{}`$ which drops slightly below unity and remains there. In this calculation we have used the proper $`T`$-dependent bridge function and the calculated $`m^{}`$ is in good agreement with the QMC-based $`m^{}`$. This might be somewhat coincidental, as the QMC results are also based on sensitive approximations. However, it means that we do have a $`B(r)`$ which is consistent with current QMC results, and hence may be used with greater confidence in studying thick-2D systems. Another point in favour of our model of $`B(r)`$ is seen in the local-field factor (LFF) of the ideal 2DES response function. A study of the LFF of the 2D responselocf shows that the formation of singlet-pair correlations is essentially complete by $`r_s5`$, and after that the structure of the fluid remains more or less unchanged, until the SPT is reached. The hard-sphere model of $`B(r)`$ provided a satisfactory description of the short-ranged features of the 2DES-LFF. The rapid rise in $`m^{}`$ up to $`r_s5`$ and the subsequent slow-down is probably related to the formation and persistence of the singlet structure in the 2D fluid revealed by the form of the LFF.
In Fig. 7 we present a comparison of various theoretical calculations of the effective mass $`m_H^{}`$ of the electrons in the HIGFET. The PT calculations of Zhang et al., and Asgari et al., show a strong decrease of $`m^{}`$ from the PT values in the ideal 2DES. Our calculations, using the โunperturbed-$`g`$โ approximation, Eq. 26, lead to a $`m_H^{}`$ which is only slightly reduced by the thickness effect. This $`m_H^{}`$ curve is in close to that of Asgari et al. This is clearly a numerical accident. According to Asgari et. al., the difference between the ideal $`m^{}`$ and $`m_H^{}`$ increase as $`r_s`$ increases. In our calculation using the โunperturbed-$`g`$โ approximation, the difference, already quite small, seems to diminish as $`r_s`$ increases. In fact, as $`r_s`$ increases, the ratio $`w/r_s`$ of the HIGFET layer decreases and the thickness effect may be expected to decrease, unless the difference between $`m^{}`$ and $`m_H^{}`$ is driven by some other effect. This โother effectโ is revealed by giving up the โunperturbed-$`g`$โ approximation, and using the full thick-layer 2DES pair-distribution function at finite $`T`$, calculated using the CHNC, to evaluate the total free energy $`F(r_s,T)`$ of the quasi-2D system, and hence the $`m_H^{}`$. In Fig. 8 we display the PDF of the quasi-2DES of a HIGFET at $`r_s=5`$, and compare it with the PDF of the ideal 2DES. The difference between the ideal and quasi systems is embodied in the form factor $`F(r)`$. The reduced Coulomb repulsion at small-$`r`$ leads to a large pile-up of electrons around the electron at the origin. This means, the electron has to drag this charge pileup and this contributes an enhanced $`m^{}`$ to the thermodynamic and transport properties of the quasi 2DES.
The experimental results of Tan et. al., for $`m^{}`$ show an increase of $``$ 150% between $`r_s`$=3 and $`r_s=6`$. Our results from the full CHNC calculation, as well as the perturbation results of Asgari et al., are shown in the lower panel of Fig. 8. Given the failure of the PT calculations to reproduce the QMC-based $`m^{}`$ for the ideal 2D, it is difficult to evaluate the reliability of the PT-based $`m_H^{}`$. The PT-overestimate of $`m^{}`$ of the ideal 2D is typical of the RPA-like character of these theories which are likely to predict spin transitions at relatively high densities. Also, we believe that if the same PT prescriptions were used to evaluate the one-top value $`g(0)`$ of the PDFS of the 2DES and the quasi-2DES, another measure of the short-comings of the PT methods would be revealed.
### V.2 Enhanced spin susceptibility and the Lande-$`g`$ factor.
The product $`m^{}g^{}`$ is given by the ratio of the static spin susceptibility $`\chi _s`$ to the ideal (Pauli) spin susceptibility $`\chi _P`$. The long wavelength limit of the static response functions are connected with the compressibility or the spin-stiffness via the second derivative of the total energy with respect to the density or the spin polarization2valley . De Palo et al.morsen have calculated $`m^{}g^{}`$ from the QMC pair distribution functions and shown that they obtain quantitative agreement with the data for very narrow 2-D systemsvakili as well as for the thicker systems found in HIGFETSzhu . The CHNC PDFs are close approximations to QMC results, and when used in Eq. 26, yield correction energies which are in good agreement with the energies obtained by De Palo et aldepalo . In Ref. 2valley we showed that the rapid enhancement of $`m^{}g^{}`$ in Si/SiO<sub>2</sub> 2DES is a consequence of the increase in $`m^{}`$ with $`r_s`$, and that the $`g^{}`$ does nอกot increase with $`r_s`$ because there is no spin-phase transition in the 2-valley case. In the HIGFET system there is a slightly delayed SPT, as seen in Fig. 5. Hence $`m^{}g^{}`$ increases with $`r_s`$, while $`m^{}`$ also increases quite rapidly, due to the enhanced โon-topโ correlations shown (Fig. 8) in the PDF of the quasi-2DES. The resulting $`g^{}`$ of the HIGFET is shown in Fig. 9.
## VI Conclusion.
We have presented a detailed study of the effect of many-body interactions in quasi-2D electron layers using a single theoretical framework which involves the calculation of the pair-distribution functions of the system via a classical representation of the quantum fluid. A procedure for replacing the inhomogeneous transverse distributions via a constant-density model, i.e., an equivalent homogeneous distribution, has also been demonstrated. Easy to use parametrized fit formulae for the exchange energy at zero and finite-$`T`$ have been presented. A simple numerical scheme for calculating the correlation energy of a thick 2D layer,, via a 3-D slab model combined with a 1-D rod model, has also been demonstrated. We find that the thickness effect on the spin-phase transition etc., provides a clear picture of the changes in the spin-susceptibility enhancement leading to a strong increase in the $`g`$-factor, while $`m^{}`$ is increased due to the enhancement of the โon-topโ correlations arising from the reduction of the Coulomb potential in thick layers. However, unlike in the case of the effective mass data for Si/SiO<sub>2</sub> systemsshash ; eplmass , these results do not provide good quantitative agreement with the effective-mass data for HIGFETS recently reported by Tan et al. This may be due to our use of the ideal 2DES bridge function even for the HIGFETS.
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# Selection Rules for Black-Hole Quantum Transitions
\[
## Abstract
We suggest that quantum transitions of black holes comply with selection rules, analogous to those of atomic spectroscopy. In order to identify such rules, we apply Bohrโs correspondence principle to the quasinormal ringing frequencies of black holes. In this context, classical ringing frequencies with an asymptotically vanishing real part $`\omega _R`$ correspond to virtual quanta, and may thus be interpreted as forbidden quantum transitions. With this motivation, we calculate the quasinormal spectrum of neutrino fields in spherically symmetric black-hole spacetimes. It is shown that $`\omega _R0`$ for these resonances, suggesting that the corresponding fermionic transitions are quantum mechanically forbidden.
\]
The necessity of a quantum theory of gravity was already recognized in the early days of quantum mechanics and general relativity. However, despite the flurry of research in this field we still lack a complete theory of quantum gravity. In many respects the black hole plays the same role in gravitation that the atom played in the nascent of quantum mechanics . It is therefore believed that black holes may play a major role in our attempts to shed light on the nature of a quantum theory of gravity.
The quantization of black holes was proposed long ago by Bekenstein , based on the remarkable observation that the horizon area of a non-extremal black hole behaves as a classical adiabatic invariant. In the spirit of the Ehrenfest principle โ any classical adiabatic invariant corresponds to a quantum entity with a discrete spectrum, and based on the idea of a minimal increase in black-hole surface area , Bekenstein conjectured that the horizon area of a quantum black hole should have a discrete spectrum of the form
$$A_n=\gamma \mathrm{}_P^2n;n=1,2,3,\mathrm{},$$
(1)
where $`\gamma `$ is a dimensionless constant, and $`\mathrm{}_P=(G\mathrm{}/c^3)^{1/2}`$ is the Planck length (we use gravitational units in which $`G=c=1`$ henceforth). This type of area quantization has since been reproduced based on various other considerations (see e.g., for a detailed list of references).
Mukhanov and Bekenstein have suggested an independent argument in order to determine the value of the coefficient $`\gamma `$. In the spirit of the Boltzmann-Einstein formula in statistical physics, they relate $`g_n\mathrm{exp}[S_{BH}(n)]`$ to the number of the black hole microstates that correspond to a particular external macro-state, where $`S_{BH}`$ is the black-hole entropy. In other words, $`g_n`$ is the degeneracy of the $`n`$th area eigenvalue. Now, the thermodynamic relation between black-hole surface area and entropy, $`S_{BH}=A/4\mathrm{}`$, can be met with the requirement that $`g_n`$ has to be an integer for every $`n`$ only when
$$\gamma =4\mathrm{ln}k,$$
(2)
where $`k`$ is some natural number.
Identifying the specific value of $`k`$ requires further input. This information may emerge by applying Bohrโs correspondence principle to the (discrete) quasinormal mode (QNM) spectrum of black holes . Gravitational waves emitted by a perturbed black hole are dominated by this โquasinormal ringingโ, damped oscillations with a discrete spectrum (see e.g., for a detailed review). At late times, all perturbations are radiated away in a manner reminiscent of the last pure dying tones of a ringing bell . These black-hole resonances are the characteristic โsoundโ of the black hole itself, depending on its parameters: mass, charge and angular momentum.
It turns out that for a Schwarzschild black hole, for a given angular harmonic index $`l`$ there exist an infinite number of (complex) quasinormal frequencies, characterizing oscillations with decreasing relaxation times (increasing imaginary part) . On the other hand, it was found numerically that the real part of the Schwarzschild gravitational resonances approaches an asymptotic constant value (we normalize $`2M=1`$, and assume a time dependence of the form $`e^{i\omega t}`$),
$$\omega _n=0.0874247\frac{i}{2}\left(n+\frac{1}{2}\right).$$
(3)
Based on Bohrโs correspondence principle, it was suggested that this asymptotic real value actually equals $`\mathrm{ln}(3)/(4\pi )`$ . An analytical proof of this equality was later given in . This was followed by a calculation of the asymptotic QNM frequencies of scalar and gravitational-electromagnetic fields in the charged Reissner-Nordstrรถm (RN) spacetime .
The emission of a quantum of frequency $`\omega `$ results in a change $`\mathrm{\Delta }M=\mathrm{}\omega _R`$ in the black-hole mass. Assuming that $`\omega `$ corresponds to the asymptotically damped limit Eq. (3) , this implies a change $`\mathrm{\Delta }A=32\pi M\mathrm{\Delta }M=4\mathrm{}\mathrm{ln}3`$ in the black hole surface area. Thus, the correspondence principle, as applied to the black-hole resonances, provides the missing link, and gives evidence in favor of the value $`k=3`$. The coefficient $`\gamma =4\mathrm{ln}3`$ is a unique value, consistent with the area-entropy thermodynamic relation, with statistical physics arguments (namely, the Boltzmann-Einstein formula), and with Bohrโs correspondence principle .
Furthermore, it was later suggested to use the black-hole QNM frequencies in order to fix the value of the Immirzi parameter in Loop Quantum Gravity, a viable approach to the quantization of general relativity . The intriguing proposals outlined above have triggered a flurry of research attempting to calculate the asymptotic ringing frequencies of various types of black holes (for a detailed list of references see, e.g., ).
The discrete black-hole mass (area) spectrum implies a discrete line emission from a quantum black hole; the frequencies of the radiation quanta emitted by the black hole will be integer multiples of the fundamental frequency $`\omega _0=\mathrm{ln}3/4\pi `$ . If true, this result indicates a distortion of Hawkingโs semiclassical spectrum . Such a modification should not be met with surprise: one should bear in mind that Hawkingโs prediction of black-hole evaporation is semiclassical in the sense that the matter fields are treated quantum mechanically, but the spacetime (and the black hole itself) are treated classically. One therefore suspects that some modifications in the character of the radiation will arise when quantum properties of the black hole itself are properly taken into account. (This state of affairs is reminiscent of atomic spectroscopy: according to the classical laws of electrodynamics, an atom should have a continuous spectrum, whereas quantum mechanics allows only a discrete line emission.)
Black-hole spectroscopy.โ In light of the preceding discussion, it is very suggestive to treat black holes as quantum objects, with a quantized surface area. The analogy with fundamental objects such as atoms raises the possibility of associating black holes with other quantum phenomena. For instance, by analogy with atomic transitions, it is natural to ask whether there are selection rules which dictate the allowed black-hole quantum transitions. Perhaps such selection rules can be inferred from the black-hole QNM spectrum, utilizing Bohrโs correspondence principle.
In order to advocate this idea, we present a possible analogy between atomic spectroscopy and an intriguing feature of the black-hole QNM spectrum. Atomic transitions are constrained by the selection rule $`\mathrm{\Delta }l=\pm 1`$ , i.e. the angular momentum $`l`$ of the atom must change when emitting a radiation quantum. Somewhat similarly, the asymptotic ringing frequency of the (rotating) Kerr black hole vanishes when the azimuthal quantum number $`m`$ of the emitted field is zero . Such a quasinormal mode corresponds to a virtual quantumโ it bears no energy. This, in turn, indicates that the corresponding quantum transition is forbidden. Hence, a Kerr black hole must change its angular momentum when emitting a field quantum, in close resemblance of the atom.
Next, we point out that for a Schwarzschild black hole, the asymptotic real value of the QNM resonances is zero for fermionic fields , as opposed to the aforementioned $`\mathrm{ln}3/4\pi `$ value found for gravitational and scalar fields. Taking cognizance of the correspondence principle, this suggests that a quantized Schwarzschild black hole cannot emit a quantum with half-integer angular momentumโ such transitions seem to be forbidden.
We conjecture that this selection-rule is more general, a genuine feature of black holes. In order to examine this hypothesis, we calculate the asymptotic ringing frequencies of a fermionic field in the RN spacetime.
The dynamics of a two-component Weyl neutrino field in the RN spacetime is governed by the Teukolsky wave equation . The black hole QNMs correspond to solutions of the wave equation with the physical boundary conditions of purely outgoing waves at spatial infinity and purely ingoing waves crossing the event horizon . Such boundary conditions single out a discrete set of resonances $`\{\omega _n\}`$. The solution to the radial Teukolsky equation may be expressed as (assuming an azimuthal dependence of the form $`e^{im\varphi }`$)
$`R_{lm}`$ $`=`$ $`e^{i\omega r}(rr_{})^{1s+i\omega +i\sigma _+}(rr_+)^{si\sigma _+}`$ (5)
$`\times \mathrm{\Sigma }_{n=0}^{\mathrm{}}d_n\left({\displaystyle \frac{rr_+}{rr_{}}}\right)^n,`$
where $`r_\pm =M\pm (M^2Q^2)^{1/2}`$ are the black hole (event and inner) horizons, $`\sigma _+\omega r_+/(r_+r_{})`$, and the field spin-weight parameter $`s`$ takes the values $`s=\pm \frac{1}{2}`$ for the neutrino field.
The sequence of expansion coefficients $`\{d_n:n=1,2,3,\mathrm{}\}`$ is determined by a recurrence relation of the form
$$\alpha _nd_{n+1}+\beta _nd_n+\gamma _nd_{n1}=0,$$
(6)
with initial conditions $`d_0=1`$ and $`\alpha _0d_1+\beta _0d_0=0`$. The quasinormal frequencies are determined by the requirement that the series in Eq. (6) is convergent, that is $`\mathrm{\Sigma }d_n`$ exists and is finite .
One finds that the physical content of the recursion coefficients becomes clear when they are expressed in terms of the Bekenstein-Hawking temperature $`T_{BH}=(r_+r_{})/A`$, where $`A=4\pi r_+^2`$ is the black-hole surface area. The recursion coefficients then obtain a surprisingly simple form,
$$\alpha _n=(n+1)(n+1s2i\beta _+\omega ),$$
(7)
$`\beta _n`$ $`=`$ $`2(n+{\displaystyle \frac{1}{2}}2i\beta _+\omega )(n+{\displaystyle \frac{1}{2}}2i\omega r_+)`$ (9)
$`s{\displaystyle \frac{1}{2}}A_{lm},`$
and
$$\gamma _n=(n2i\omega )(n+s2i\beta _+\omega ),$$
(10)
where $`\beta _+(4\pi T_{BH})^1`$ is the black-hole inverse temperature. The angular separation constants are given by $`A_{lm}=l(l+1)s(s+1)`$, where $`l\mathrm{max}(|m|,|s|)`$ is the angular momentum of the field . We calculate numerically the quasinormal spectrum of neutrino fields by solving Eqs. (6)-(10) using the method of continued fractions .
We find that the neutrino quasinormal frequencies of the RN black hole exhibit a damped periodic behavior, where asymptotically $`\omega _R0`$ as $`|\omega _I|\mathrm{}`$. Figure 1 demonstrates this damped periodic dependence of $`\omega _R`$ on $`\omega _I`$, in the complex frequency plane. In figure 2 we depict the envelope of these damped oscillations, which suggests that $`\omega _R\omega _I^{1/2}`$ in the asymptotic large damping limit. Our results thus confirm that the asymptotically vanishing resonances are a general feature of neutrino fields in spherically symmetric black-hole spacetimes. According to the correspondence principle, they describe virtual, zero energy quanta, suggesting that the corresponding quantum transitions are forbidden.
Summary.โ In the spirit of Bohrโs correspondence principle, we suggest and demonstrate that asymptotically vanishing QNM frequencies may be interpreted as zero energy, forbidden quantum transitions. Motivated by this idea, we calculate the ringing frequencies of neutrino fields in the RN spacetime. It is shown that the spectrum of these black-hole resonances is characterized by a vanishing asymptotic real value. Our results raise the possibility that quantized spherically symmetric black holes cannot emit such fermionic quanta, possessing half-integer values of the angular momentum. Quantum black-hole selection rules, such as those demonstrated in this study, may provide an important clue in unveiling the underlying principles of the elusive theory of quantum gravity.
ACKNOWLEDGMENTS
The research of SH was supported by G.I.F. Foundation.
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# Algebraic Closed Geodesics on a Triaxial Ellipsoid 11footnote 1AMS Subject Classification 14H52, 37J45, 53C22, 58E10
## 1 Introduction
One of the best known classical integrable systems is the geodesic motion on a triaxial ellipsoid $`Q^3`$. By introducing ellipsoidal coordinates on $`Q`$, the problem was reduced to hyperelliptic quadratures by Jacobi () and was integrated in terms of theta-functions of a genus 2 hyperelliptic curve $`\mathrm{\Gamma }`$ by Weierstrass in .
A generic geodesic on $`Q`$ is known to be quasiperiodic and oscillate between 2 symmetric curvature lines (caustics).
It is of a certain interest to find conditions for a geodesic on $`Q`$ to be periodic (closed) and to describe such geodesics explicitly. At first sight this problem has a standard solution: by introducing actionโangle variables $`\{I_1,I_2,\varphi _1,\varphi _2\}`$ one can define frequencies
$$\dot{\varphi }_1=\mathrm{\Omega }_1(I_1,I_2),\dot{\varphi }_2=\mathrm{\Omega }_2(I_1,I_2).$$
Then the geodesic is closed if and only if the rotation number $`\mathrm{\Omega }_2/\mathrm{\Omega }_1`$ is rational.
However, the frequencies $`\mathrm{\Omega }_j`$ are known to be linear combinations of Abelian integrals on the hyperelliptic curve, hence the condition on the rotation number implies a transcendental equation on the constants of motion and the parameters of the problem. In practice, this appears to be useless for an exact description of closed geodesics.
Such a description can be made much more explicit when the hyperelliptic curve $`\mathrm{\Gamma }`$ turns out to be a covering of an elliptic curve $``$ and a certain holomorphic differential reduces to a holomorphic differential on $``$. Then the corresponding geodesic itself is a spatial elliptic curve, which covers $``$<sup>2</sup><sup>2</sup>2More precisely, it is a connected component of a real part of an elliptic curve, or a rational curve and, as an algebraic subvariety in $`^3`$ ($`^3`$), it can be represented as a connected component of the intersection of the ellipsoid $`Q`$ with an algebraic surface. In the sequel, such class of geodesics will be referred to as algebraic closed geodesics.
Conversely, one can also show that any algebraic closed geodesic on an ellipsoid must be a connected component of an elliptic or a rational curve.
Studying closed geodesics on quadrics is a classical problem. Surprisingly, we did not find any reference to its explicit solution in the classical or modern literature. Here we can only quote the paper , which studied the case of a 2-fold covering of an elliptic curve, when the solution is expressed in terms of two elliptic functions of time with different period lattices. Since the periods are generally incommensurable, the corresponding geodesics are not periodic but quasi-periodic.
Note that in the problem of periodic orbits of the Birkhoff billiard inside an ellipsoid much more progress has been made (see ).
#### Contents of the paper.
The paper proposes a simple approach to explicit description of algebraic surfaces $`๐ฑ`$ in $`^3`$ that cut out closed geodesics on $`Q`$. It is based on elements of the WeierstrassโPoncarรฉ theory of reduction of Abelian functions (see, e.g.,), addition law for elliptic functions, and the remarkable relation between geodesics on a quadric and stationary solutions of the KdV equation (the MoserโTrubowitz isomorphism) described in and recently revisited in in connection with periodic orbits of geodesic billiards on an ellipsoid.
Namely, for each genus 2 hyperelliptic tangential cover of an elliptic curve we construct a one-parameter family of plane algebraic curves (so called polhodes). Appropriate connected components of the polhodes have a form of Lissajou curves and describe closed geodesics in terms of the two ellipsoidal coordinates $`\lambda _1,\lambda _2`$ on $`Q`$. The geodesics of one and the same family are tangent to the same caustic on $`Q`$ and the parameters of the ellipsoid are functions of the moduli of the elliptic curve.
Since equations of the polhodes depend only on the symmetric functions $`\lambda _1+\lambda _2`$, $`\lambda _1\lambda _2`$, they can be rewritten in terms of Cartesian coordinates in $`^3`$ thus giving the above mentioned algebraic surfaces $`๐ฑ`$.
Our family of closed geodesic contains special ones, which have mirror symmetry with respect to principal coordinate planes in $`^3`$. In particular, for the case of 3-fold and 4-fold coverings of an elliptic curve, the special geodesics are cut out by quadratic and, respectively, cubic surfaces in $`^3`$, as illustrated in Figures 5.4 and 6.3.
Depending on how to assign the parameters of the ellipsoid and of the caustic to the branch points of the hyperelliptic curve, one can obtain closed geodesics with or without self-intersections on $`Q`$.
## 2 Linearization of the geodesic flow on the <br>ellipsoid and some elliptic solutions
We first briefly recall the integration of the geodesic motion on an $`n`$-dimensional ellipsoid
$$Q=\left\{\frac{X_1^2}{a_1}+\mathrm{}+\frac{X_{n+1}^2}{a_{n+1}}=1\right\}^{n+1}=(X_1,X_2,\mathrm{},X_{n+1}),$$
$$0<a_1<\mathrm{}<a_n<a_{n+1}.$$
Let $`t`$ be the natural parameter of the geodesic and $`\lambda _1,\mathrm{},\lambda _n`$ be the ellipsoidal coordinates on $`Q`$ defined by the formulas
$$X_i^2=a_i\frac{(a_i\lambda _1)\mathrm{}(a_i\lambda _n)}{_{ji}(a_ia_j)},i=1,\mathrm{},n+1.$$
(2.1)
In these coordinates and their derivatives $`d\lambda _k/dt`$ the total energy $`{\displaystyle \frac{1}{2}}(\dot{X},\dot{X})`$ takes a Stรคckel form, and after time re-parameterization <sup>3</sup><sup>3</sup>3For $`n=2`$ this re-parameterization was made by Weierstrass .
$$dt=\lambda _1\mathrm{}\lambda _nds,$$
(2.2)
the evolution of $`\lambda _i`$ is described by quadratures
$$\frac{\lambda _1^{k1}d\lambda _1}{2\sqrt{R(\lambda _1)}}+\mathrm{}+\frac{\lambda _n^{k1}d\lambda _n}{2\sqrt{R(\lambda _n)}}=\{\begin{array}{cc}\hfill ds\text{ for }& k=1,\hfill \\ \hfill 0\text{ for }& k=2,\mathrm{},n,\hfill \end{array}$$
(2.3)
$$R(\lambda )=\lambda (\lambda a_1)\mathrm{}(\lambda a_{n+1})(\lambda c_1)\mathrm{}(\lambda c_n),$$
where $`c_k`$ are constants of motion.
This implies integrability of the system by the Liouville theorem. The generic invariant varieties of the flow are $`n`$-dimensional tori with a quasiperiodic motion. The corresponding geodesics are tangent to one and the same set of $`n1`$ confocal quadrics $`Q_{c_1},\mathrm{},Q_{c_{n1}}`$ of the confocal family
$$Q_c=\left\{\frac{X_1^2}{a_1c}+\mathrm{}+\frac{X_{n+1}^2}{a_{n+1}c}=1\right\}^{n+1}.$$
(2.4)
In particular, generic geodesics on a 2-dimensional ellipsoid fill a ring bounded by caustics, the lines of intersection of $`Q`$ with confocal hyperboloid $`Q_{c_1}`$.
The quadratures (2.3) involve $`n`$ independent holomorphic differentials on the genus $`n`$ hyperelliptic curve $`\mathrm{\Gamma }=\{\mu ^2=R(\lambda )\}`$,
$$\omega _k=\frac{\lambda ^{k1}d\lambda }{2\sqrt{R(\lambda )}},k=1,\mathrm{},n$$
(2.5)
and give rise to the AbelโJacobi map of the $`n`$-th symmetric product $`\mathrm{\Gamma }^{(n)}`$ to the Jacobian variety of $`\mathrm{\Gamma }`$,
$$\underset{P_0}{\overset{P_1}{}}\omega _k+\mathrm{}+\underset{P_0}{\overset{P_n}{}}\omega _k=u_k,P_k=(\lambda _k,\sqrt{R(\lambda _k)})\mathrm{\Gamma },$$
(2.6)
where $`u_1,\mathrm{},u_n`$ are coordinates on the universal covering of Jac$`(\mathrm{\Gamma })`$ and $`P_0`$ is a fixed basepoint, which we choose to be the infinity point $`\mathrm{}`$ on $`\mathrm{\Gamma }`$.
Since $`u_1=s+`$const and $`u_2,\mathrm{},u_n=`$const, the geodesic motion in the new parameterization is linearized on the Jacobian variety of $`\mathrm{\Gamma }`$.
The inversion of the map (2.6) applied to formulas (2.1) leads to the following parameterization of a generic geodesic in terms of $`n`$-dimensional theta-functions $`\theta (w_1,\mathrm{},w_n)`$ associated to the curve $`\mathrm{\Gamma }`$,
$$X_i(s)=\varkappa _i\frac{\theta [\eta _i](w_1,\mathrm{},w_n)}{\theta [\mathrm{\Delta }](w_1,\mathrm{},w_n)},i=1,\mathrm{},n+1,$$
(2.7)
where $`\mathrm{\Delta }=(\delta ^{\prime \prime },\delta ^{})`$, $`\eta _i=(\eta _i^{\prime \prime },\eta _i^{})^{2n}/2^{2n}`$ are certain half-integer theta-characteristics, the arguments $`w_1,\mathrm{},w_n`$ depend linearly on $`u_1,\mathrm{},u_n`$, and therefore on $`s`$, and $`\varkappa _i`$ are constant factors depending on the moduli of $`\mathrm{\Gamma }`$ only.
For the classical Jacobi problem $`(n=2)`$, the complete theta-functional solution was presented in , and, for arbitrary dimensions, in , whereas a complete classification of real geodesics on $`Q`$ was made in .
#### Periodicity problem and a solution in terms of elliptic functions.
As mentioned in Introduction, we restrict ourselves with the case when a geodesic is periodic in the complex parameter $`s`$, namely, double-periodic. This implies that the solution (2.7) can be expressed in terms of elliptic functions of $`s`$.
As an example, following von Braunmuhl , consider the geodesic problem on 2-dimensional quadric ($`n=2`$) and suppose that the parameters $`a_i,c_j`$ in (2.3) are such that the curve $`\mathrm{\Gamma }`$ becomes birationally equivalent to the following canonical curve
$$w^2=z(z1)(z\alpha )(z\beta )(z\alpha \beta ),$$
$`\alpha ,\beta `$ being arbitrary positive constant. Then, as widely described in the literature (see, e.g., ), $`\mathrm{\Gamma }`$ covers two different elliptic curves
$$\mathrm{}_\pm =\{W_\pm ^2=Z_\pm (1Z_\pm )(1k_\pm Z_\pm )\},k_\pm ^2=\frac{(\sqrt{\alpha }\sqrt{\beta })^2}{(1\alpha )(1\beta )}$$
with covering relations
$`Z_+=Z_{}`$ $`={\displaystyle \frac{(1\alpha )(1\beta )z}{(z\alpha )(z\beta )}},`$ (2.8)
$`W_\pm `$ $`=\sqrt{(1\alpha )(1\beta )}{\displaystyle \frac{z\sqrt{\alpha \beta }}{(z\alpha )^2(z\beta )^2}}w.`$
Thus, $`\mathrm{\Gamma }`$ is a 2-fold covering of $`\mathrm{}_{}`$ and $`\mathrm{}_+`$.
Both holomorphic differentials $`\omega _1,\omega _2`$ on $`\mathrm{\Gamma }`$ reduce to linear combinations of the holomorphic differentials on $`\mathrm{}_+`$ and $`\mathrm{}_{}`$, namely
$$\frac{dZ_\pm }{W_\pm }=\frac{z\sqrt{\alpha \beta }}{w}dz.$$
Then a linear combination of equations (2.3) for $`n=2`$ yields
$`{\displaystyle _{\mathrm{}}^{\lambda _1}}{\displaystyle \frac{z\sqrt{\alpha \beta }}{w}}๐z+{\displaystyle _{\mathrm{}}^{\lambda _2}}{\displaystyle \frac{z\sqrt{\alpha \beta }}{w}}๐z`$ $`=s\sqrt{\alpha \beta }+\text{const},`$
$`{\displaystyle _{\mathrm{}}^{\lambda _1}}{\displaystyle \frac{z+\sqrt{\alpha \beta }}{w}}๐z+{\displaystyle _{\mathrm{}}^{\lambda _2}}{\displaystyle \frac{z+\sqrt{\alpha \beta }}{w}}๐z`$ $`=s\sqrt{\alpha \beta }+\text{const}.`$
Inversion of these quadratures lead to solutions for $`X_i`$ in terms of elliptic functions of the curves $`\mathrm{}_\pm `$, whose arguments both depend on the time parameter $`s`$. Then, since their periods are generally incommensurable, the corresponding geodesics remain to be quasi-periodic.
This observation shows that not any case of covering $`\mathrm{\Gamma }`$ to an elliptic curve results in closed geodesics on $`Q`$. In the next section we consider other types of coverings and obtain sufficient condition for a geodesic to be an elliptic curve.
## 3 Hyperelliptic tangential covers and closed <br>geodesics on an $`n`$-dimensional ellipsoid
Consider a genus $`n`$ compact smooth hyperelliptic surface $`G`$, whose affine part $`G_A^2=(z,w)`$ is given by equation
$$w^2=R_{2n+1}(z),$$
$`R_{2n+1}(z)`$ being a polynomial of degree $`2n+1`$. The curve $`G`$ is obtained from $`G_A`$ by gluing the infinite point $`\mathrm{}`$. Let $`\{\mathrm{\Omega }_1(P),\mathrm{},\mathrm{\Omega }_n(P)\}`$, $`P=(z,w)G`$ be a basis of independent holomorphic differentials on $`G`$. One can also write $`\mathrm{\Omega }_j(P)=\varphi _j(P)d\tau `$, where $`\tau `$ is a local coordinate in a neibourhood of $`P`$. Next, let $`\mathrm{\Lambda }`$ be the lattice in $`^n`$ generated by $`2n`$ independent period vectors $`(\mathrm{\Omega }_1,\mathrm{},\mathrm{\Omega }_n)^T`$.
The curve $`G`$ admits a canonical embedding into its Jacobian variety Jac$`(G)=^n(u_1,\mathrm{},u_n)/\mathrm{\Lambda }`$,
$$P๐(P)=_{\mathrm{}}^P(\mathrm{\Omega }_1,\mathrm{},\mathrm{\Omega }_n)^T,$$
(3.1)
so that $`\mathrm{}`$ is mapped into the neutral point (origin) in Jac$`(G)`$ and
$$๐=\frac{d}{d\tau }๐(๐ซ)|_{P=\mathrm{}}=(\varphi _1(\mathrm{}),\mathrm{},\varphi _n(\mathrm{}))^T$$
is the tangent vector of $`G`$Jac$`(G)`$ at the origin.
Now assume that $`G`$ is an $`N`$-fold covering of an elliptic curve $``$, which we represent in the canonical Weierstrass form
$$=\left\{(\mathrm{}^{}(u))^2=4\mathrm{}^3(u)g_2\mathrm{}(u)g_34(\mathrm{}e_1)(\mathrm{}e_2)(\mathrm{}e_3)\right\}.$$
(3.2)
Here $`\mathrm{}(u)=\mathrm{}(u\omega ,\omega ^{})`$ denotes the Weierstrass elliptic function with half-periods $`\omega ,\omega ^{}`$ and $`u/\{\omega +\omega ^{}\}`$. The parameters $`g_2,g_3`$ provide moduli of the curve.
Assume also that under the covering map $`\pi :G`$ the infinite point $`\mathrm{}G`$ is mapped to $`u=0`$.
In the sequel we concentrate on hyperelliptic tangential coverings $`\pi :G`$, when $``$ admits the following canonical embedding onto Jac$`(G)`$
$$u(u_1,\mathrm{},u_n)=u๐.$$
That is, the images of $``$ and $`G`$ in Jac$`(G)`$ are tangent at the origin <sup>4</sup><sup>4</sup>4As follows from this definition, the 2-fold covers (2.8) are not hyperellipticallly tangential.. The motion of tangential covering was introduced in in connection with elliptic solutions of the KdV equation (see also ).
Namely, let $`\theta (u_1,\mathrm{},u_n)`$ be the theta-function associated to the covering curve $`G`$ and $`\mathrm{\Theta }`$ be the theta-divisor, codimension one subvariety of $`\mathrm{Jac}(\mathrm{\Gamma })`$ defined by equation $`\theta [\mathrm{\Delta }](๐ฎ)=0`$, where $`[\mathrm{\Delta }]`$ is the special theta-characteristic in the solution (2.7).
###### Theorem 3.1.
() For an arbitrary vector $`๐^g`$, the transcendental equation
$$\theta (๐x+๐)=0,x,$$
(3.3)
has exactly $`N`$ solutions $`x=q_1(๐),\mathrm{},x=q_N(๐)`$ (possibly, with multiplicity).
That is, the complex flow on Jac$`(G)`$ in $`๐`$-direction intersects the theta-divisor $`\mathrm{\Theta }`$ or any its translate at a finite number of points. This property is exceptional: for a generic hyperelliptic curve $`G`$ the number of such intersections is infinite.
Note that in the local coordinates $`u_1,\mathrm{},u_n`$ on Jac($`G`$) corresponding to the standard basis of holomorphic differentials
$$\overline{\omega }_k=\frac{z^{k1}dz}{w},k=1,\mathrm{},n,$$
(3.4)
one has $`๐=(0,\mathrm{},0,2)^T`$.
According to the Poincarรฉ reducibility theorem (see e.g., ), apart from the curve $``$, the Jacobian of $`G`$ contains an $`(n1)`$-dimensional Abelian subvariety $`๐_{n1}`$. For $`n=2`$ the subvariety is just another elliptic curve covered by $`G`$.
Notice that for the case $`n=2`$, explicit algebraic expressions of the covers and coefficients of hyperelliptic curves are known for $`N8`$ (see ).
#### Double periodic geodesics on an ellipsoid.
The algebraic geometrical property described by Theorem 3.1 gives a tool for a description of double-periodic geodesic flow on the $`n`$-dimensional quadric $`Q`$, which is linearized on the Jacobian of the hyperelliptic curve $`\mathrm{\Gamma }`$ in Section 2. Namely, let the genus $`n`$ curves $`G`$ and $`\mathrm{\Gamma }`$ are related via birational transformation of the form
$$\lambda =\frac{\alpha }{(z\beta )},\mu =\frac{w}{(z\beta )^{n+1}},$$
(3.5)
where $`(\beta ,0)`$ is a finite Weierstrass point on $`G`$ and $`\alpha `$ is an arbitrary positive constant. Then the following theorem proved in holds.
###### Theorem 3.2.
To any hyperelliptic tangential cover $`G`$ such that all the Weierstrass points of $`G`$ are real, one can associate an $`(n1)`$-parametric family of different closed real geodesics on an $`n`$-dimensional ellipsoid $`Q`$ that are tangent to the same set of confocal quadrics $`Q_{c_1},\mathrm{},Q_{c_{n1}}`$. The parameters of the ellipsoid ($`a_i`$) and of the quadrics ($`c_j`$) are related to branch points of $`G`$ via the transformation (3.5).
#### Remark.
It is natural to consider a closed geodesic as a curve on $`Q`$ and not as a periodic solution of the geodesic equations that depends on the initial point on the curve as on a parameter. That is, we disregard this parameter in the above family of closed real geodesics.
Proof of Theorem 3.2. The transformation (3.5) sends the points $`\mathrm{}`$ and $`(\beta ,0)`$ on $`G`$ to the Weierstrass points $`๐ช=(0,0)`$ and, respectively, $`\mathrm{}`$ on $`\mathrm{\Gamma }`$. Then, identifying the curves $`G`$ and $`\mathrm{\Gamma }`$, as well as their Jacobians, we find that the $`๐`$-flow on Jac($`G`$), which is tangent to the canonically embedded hyperelliptic curve $`\mathrm{\Gamma }\text{Jac}(\mathrm{\Gamma })`$ at $`\mathrm{}`$, is represented as the flow on Jac($`\mathrm{\Gamma }`$) which is tangent to the embedded $`\mathrm{\Gamma }\mathrm{Jac}(\mathrm{\Gamma })`$ at $`๐ช`$, and vice versa. In the coordinates on Jac$`(\mathrm{\Gamma })`$ corresponding to the basis (2.5), the latter flow has direction $`(1,0,\mathrm{},0)^T`$ and thus coincides with the linearized geodesic flow on $`Q`$.
This remarkable relation was first described in as the MoserโTrubowitz isomorphism between stationary $`n`$-gap solutions of the KdV equation and generic (quasiperiodic) geodesics on an $`n`$-dimensional quadric.
Next, let us fix a real constant $`d`$ and the confocal quadric $`Q_d`$ of the family (2.4) such that the geodesics with the constants of motion $`c_1,\mathrm{},c_{n1}`$ have a non-empty intersection with $`QQ_d`$. In view of (2.1), when a geodesic $`X(s)`$ intersects $`QQ_d`$, one of the points $`P_i=(\lambda _i,\mu _i)`$ on the curve $`\mathrm{\Gamma }`$ (without loss of generality we choose it to be $`P_n`$) coincides with one of the points $`E_{d\pm }=(d,\pm \sqrt{R(d)})`$. Under the AbelโJacobi map (2.6) with $`P_0=\mathrm{}`$, the condition $`P_n=E_{d\pm }`$ defines two translates of the theta-divisor
$$\mathrm{\Theta }_{d\pm }=\{\theta [\mathrm{\Delta }](๐ฎq/2)=0\}\mathrm{Jac}(\mathrm{\Gamma }),q=_{E_d}^{E_{d+}}(\mathrm{\Omega }_1,\mathrm{},\mathrm{\Omega }_n)^T^n.$$
A geodesic is doubly-periodic if and only if it intersects $`QQ_d`$ at a finite number of complex points. In this case the linearized flow on $`\mathrm{Jac}(\mathrm{\Gamma })`$ must intersect $`\mathrm{\Theta }_{d\pm }`$ at a finite set of points too. In view of the MoserโTrubowitz isomorphism and Theorem 3.1, this holds if $`\mathrm{\Gamma }`$ is a hyperelliptic tangential cover of an elliptic curve $``$. Then, under the transformation (3.5) with an appropriate $`\beta `$, the real Weierstrass points on $`\mathrm{\Gamma }`$ give real and positive parameters $`a_i,c_j`$ of the doubly-periodic geodesic.
Finally, there is an $`(n1)`$-dimensional family of elliptic curves $``$ in Jac($`G`$), which is locally parameterized by points of their intersection with the Abelian subvariety $`๐_{n1}`$. This gives rise to an $`(n1)`$-dimensional family of the doubly-periodic geodesics. $`\overline{)}`$
#### Remark.
Since for any chosen $`N`$-fold tangential cover $`G`$ the branch points of $`G`$ are functions of the two moduli $`g_2,g_3`$, the parameters $`a_i,c_j`$ are uniquely determined by them and by the rescaling factor $`\alpha `$ in (3.5). This implies that not any ellipsoid $`Q`$ may have doubly-periodic geodesics associated with the given degree of covering as described by Theorem 3.2. One can show that even in the simplest case of a triaxial ellipsoid ($`n=2`$) and $`N=3`$ or 4, for any fixed positive $`a_1,a_2`$ there exists only a finite number of possible $`c,a_3`$ for which the geodesics are doubly-periodic<sup>5</sup><sup>5</sup>5Explicit algebraic conditions on $`a_1,a_2,a_3,c`$ for the case of 3- and 4-fold tangential covers were
presented in ..
Naturally, this does not exclude the existence of such geodesics for other degrees of tangential coverings or those obtained from a periodic flow on Jac($`G`$) via a birational transformation different from (3.5), or even just closed geodesics, which are not doubly-periodic. However, the latter, if exist, cannot be algebraic curves in view of the following property.
###### Lemma 3.3.
Any algebraic closed geodesic on ellipsoid $`Q^n`$ is a connected component of an elliptic or rational curve.
Proof. Let a closed geodesic be a connected component of an algebraic curve $`๐`$. Since the geodesic flow on $`Q`$ is linearized on an unramified covering of Jac$`(\mathrm{\Gamma })`$, $`๐`$ must be an unramified covering of an algebraic curve $`๐_0\mathrm{Jac}(\mathrm{\Gamma })`$ and and, moreover, $`๐_0`$ must be a one-dimensional Abelian subvariety. Then, if $`\mathrm{\Gamma }`$ is a regular curve and, therefore, Jac$`(\mathrm{\Gamma })`$ is compact, $`๐_0`$ can be only elliptic. If $`\mathrm{\Gamma }`$ has singularities (when, for example, $`c_j=a_i`$ and the geodesic lies completely in hyperplane $`X_i=0`$) and its generalized Jacobian is not compact, then $`๐_0`$ can be also a rational curve. In both cases $`๐`$, as an unramified covering of $`๐_0`$, can be only elliptic or a reducible rational curve. $`\overline{)}`$
In the case $`n=2`$ the algebraic closed geodesics on a triaxial ellipsoid can explicitly be expressed in terms of symmetric functions of the two ellipsoidal coordinates $`\lambda _1,\lambda _2`$ on $`Q`$. As a result, such geodesics can be rewritten in terms of Cartesian coordinates in $`^3`$. We shall describe this procedure in the next section.
## 4 Genus 2 hyperelliptic tangential covers, algebraic polhodes, and cutting algebraic surfaces in $`^3`$
Suppose that the genus 2 hyperelliptic curve $`G`$
$$w^2=\underset{k=1}{\overset{5}{}}(zb_k)$$
(4.1)
is an $`N`$-fold tangential covering of the elliptic curve $``$ in (3.2). Then, according to the Poincarรฉ reducibility theorem, $`G`$ is also an $`N`$-fold covering of another elliptic curve
$$_2=\left\{W^2=\left(4Z^3G_2ZG_3\right)4(ZE_1)(ZE_2)(ZE_3)\right\},$$
the parameters $`G_2,G_3`$ being functions of the moduli $`g_2,g_3`$.
Let $`U`$ be uniformization parameter such that $`Z=\stackrel{~}{\mathrm{}}(U)`$, $`W=\stackrel{~}{\mathrm{}}^{}(U)`$, and $`\stackrel{~}{\mathrm{}}`$ is the Weierstrass function associated to the curve $`_2`$. As above, assume that the point $`\mathrm{}G`$ is mapped to $`U=0`$. Then one can show that the map $`\pi :G_2`$ is described by formulas
$$Z=๐ต(z),W=w^k๐ฒ(z),$$
(4.2)
where $`k`$ is a positive odd integer number and $`๐ต(z),๐ฒ(z)`$ are rational functions of $`z`$ such that $`๐ต(\mathrm{})=\mathrm{}`$. The second relation in (4.2) implies that the Weierstrass points on $`G`$ are mapped to branch points on $`_2`$.
Consider the canonical embedding of $`G`$ to its Jacobian variety $`^2=(u_1,u_2)/\mathrm{\Lambda }`$,
$$P=(z,w)๐(z,w)=(_{\mathrm{}}^P\frac{dz}{w},_{\mathrm{}}^P\frac{zdz}{w})^T.$$
The image of the embedding is the theta-divisor $`\mathrm{\Theta }`$ that passes through the origin in Jac$`(G)`$ and is tangent to vector $`๐=(0,1)^T`$.
The second covering $`\pi :G_2`$ is lifted to the Jacobian variety of $`G`$. Namely, for any point $`๐ฌ_2`$ and $`\pi ^1(๐ฌ)=\{P^{(1)},\mathrm{},P^{(N)}\}G`$, one has
$$\underset{\mathrm{}}{\overset{P^{(i)}}{}}\frac{dz}{w}=\kappa \underset{\mathrm{}}{\overset{๐ฌ}{}}\frac{dZ}{W},$$
(4.3)
where $`\kappa `$ is a constant rational number depending on the degree $`N`$ only. This implies that $`z`$-coordinates of the $`N`$ points of intersection of a complex $`u_2`$-line ($`u_1=`$const) with $`G=\mathrm{\Theta }\text{Jac}(G)`$ are the roots of the first equation in (4.2) with $`Z=Z(๐ฌ)`$ (see also ).
Now let $`P_1=(z_1,w_1),P_2=(z_1,w_2)G`$ and consider the full AbelโJacobi map
$$๐(P_1)+๐(P_2)=(u_1,u_2)^T.$$
(4.4)
Assume that $`u_1,u_2`$ evolve according to $`๐`$-flow, that is $`u_1=`$const. Hence $`z_1,z_2`$ satisfy the equations
$$\dot{z}_1=\frac{w_1}{z_1z_2},\dot{z}_2=\frac{w_1}{z_2z_1}.$$
(4.5)
This imposes a relation between coordinates of $`P_1`$ and $`P_2`$ on $`G`$. In the generic case, the relation is transcendental one and the coordinates are quasiperiodic functions of time. However, if $`G`$ is a tangential covering of an elliptic curve, then the relation becomes algebraic and can be found explicitly in each case of covering. Namely, let us set
$$U_1=\underset{\mathrm{}}{\overset{\pi (P_1)}{}}\frac{dZ}{W},U_2=\underset{\mathrm{}}{\overset{\pi (P_2)}{}}\frac{dZ}{W},\text{and}U_{}=U_1+U_2.$$
(4.6)
In view of (4.3) and the condition $`u_1=`$const, the first equation in (4.4) implies $`U_{}=`$const.
Next, due to the addition theorem for elliptic functions,
$$\left|\begin{array}{ccc}1& \stackrel{~}{\mathrm{}}(U_1)& \stackrel{~}{\mathrm{}}^{}(U_1)\\ 1& \stackrel{~}{\mathrm{}}(U_2)& \stackrel{~}{\mathrm{}}^{}(U_2)\\ 1& \stackrel{~}{\mathrm{}}(U_{})& \stackrel{~}{\mathrm{}}^{}(U_{})\end{array}\right|=0,$$
(4.7)
or, in the integral form
$$\stackrel{~}{\mathrm{}}(U_{})+\stackrel{~}{\mathrm{}}(U_1)+\stackrel{~}{\mathrm{}}(U_2)=\frac{1}{4}\left[\frac{\stackrel{~}{\mathrm{}}^{}(U_1)\stackrel{~}{\mathrm{}}^{}(U_2)}{\stackrel{~}{\mathrm{}}(U_1)\stackrel{~}{\mathrm{}}(U_2)}\right]^2,$$
the coordinates $`Z_1=\stackrel{~}{\mathrm{}}(U_1),Z_2=\stackrel{~}{\mathrm{}}(U_2)`$ are subject to the constraint
$`2G_3+G_2\left(Z_1+Z_2\right)`$ $`+4\mathrm{}(U_{})\left(Z_2Z_1\right)^24Z_2Z_1\left(Z_1+Z_2\right)`$
$`=2\sqrt{4Z_1^3G_2Z_1G_3}\sqrt{4Z_2^3G_2Z_2G_3}.`$
Then, taking square of both sides, simplifying, factoring out $`(Z_1Z_2)^2`$, and replacing $`Z_1,Z_2`$ by the expressions $`๐ต(z_1),๐ต(z_2)`$ from (4.2), we arrive at generating equation
$`16`$ $`\left(๐ต(z_2)๐ต(z_1)\right)^2\stackrel{~}{\mathrm{}}_{}^2`$
$`+\left[16G_3+8G_2\left(๐ต(z_1)+๐ต(z_2)\right)32๐ต(z_1)๐ต(z_2)(๐ต(z_1)+๐ต(z_2))\right]\stackrel{~}{\mathrm{}}_{}`$
$`+16G_3\left(๐ต(z_1)+๐ต(z_2)\right)+8G_2๐ต(z_1)๐ต(z_2)+16๐ต(z_1)^2๐ต(z_2)^2+G_2^2=0.`$ (4.8)
Written in terms of $`Z_1,Z_2`$, it defines an elliptic curve isomorphic to $`_2`$ for any $`\stackrel{~}{\mathrm{}}_{}`$.
In terms of $`z_1,z_2`$, the generating equation gives a family of algebraic curves $`_\stackrel{~}{\mathrm{}}_{}(z_1,z_2)`$, which we call polhodes. They are symmetric with respect to the diagonal $`z_1=z_2`$, as expected, and, for a generic $`\stackrel{~}{\mathrm{}}_{}`$, has degree $`4N`$<sup>6</sup><sup>6</sup>6 However, as seen from the structure of (4.8), a fixed generic $`z_1`$ (and $`u_1`$) results in $`2N`$ (complex) solutions $`z_2`$.. A polhode describes an algebraic relation between $`z`$-coordinates of the divisor $`P_1=(z_1,w_1),P_2=(z_1,w_2)`$ on $`G`$, which holds under the $`u_2`$-flow on Jac$`(G)`$. The parameter $`\stackrel{~}{\mathrm{}}_{}=\stackrel{~}{\mathrm{}}(U_{})`$ plays the role of a constant phase of the flow.
The polhodes thus can be regarded as ramified coverings of $`_2`$ and, therefore, in general, have genus $`>1`$.
#### Real finite asymmetric part of polhodes.
Suppose that all the roots of the degree 5 polynomial in (4.1) are real and set
$$b_1<b_2<\mathrm{}<b_5.$$
(4.9)
Assume that the variables $`z_1,z_2`$ range in finite segments $`[b_i,b_j]`$, where both $`w_1,w_2`$ are real and finite. Taking in mind applications to problems of dynamics, we also assume that these segments are different and $`z_1<z_2`$. Then the motion of the point $`(z_1,z_2)`$ is bounded in the unique square domain
$$S=\{b_2z_1b_3,b_4z_2b_5\}.$$
Let also $`\stackrel{~}{\mathrm{}}_{}`$ in (4.8) be real. The part of polhode $`_\stackrel{~}{\mathrm{}}_{}(z_1,z_2)`$ that lies in $`S`$ will be called the real asymmetric part of $`_\stackrel{~}{\mathrm{}}_{}`$. At the vertices of the domain both $`w_1,w_2`$ equal zero. Then, in view of equations (4.5), this part of the polhode is tangent to the sides of $`S`$ or passes through some of its vertices.
###### Lemma 4.1.
If $`U_{}`$ is such that in the domain $`S`$ one of the following relations holds
$$\stackrel{~}{\mathrm{}}(U_{})=๐ต(z_1),\text{or}\stackrel{~}{\mathrm{}}(U_{})=๐ต(z_2),$$
(4.10)
then the real asymmetric part of $`_\stackrel{~}{\mathrm{}}_{}`$ is empty.
Proof. Indeed, in view of (4.6), condition $`\stackrel{~}{\mathrm{}}_{}=Z(z_1)`$ implies $`U_2=0`$. Hence, for the above value of $`z_1[b_2,b_3]`$, the coordinate $`z_2`$ must be infinite. If the component of $`_\stackrel{~}{\mathrm{}}_{}`$ in $`S`$ is not empty, then the polhode must intersect the boundary of $`S`$, which is not possible. Hence this component is empty.
If the second condition in (4.10) is satisfied, the proof goes along similar lines. $`\overline{)}`$
In view of the above lemma, we also assume that the constant parameter $`\stackrel{~}{\mathrm{}}_{}`$ lies in a segment on $``$ where neither of the conditions (4.10) is satisfied, which is one of the gaps $`[E_\alpha ,E_\beta ]`$, $`[\mathrm{},E_1]`$, $`[E_3,\mathrm{}]`$.
#### Special polhodes.
If the parameter $`\stackrel{~}{\mathrm{}}_{}`$ in (4.8) coincides with a branch point of $`_2`$, then the equation of the polhode simplifies.
In the first obvious case $`\stackrel{~}{\mathrm{}}_{}=\mathrm{}`$ the generating equation (4.8) reduces to $`Z(z_1)Z(z_2)=0`$. Since $`Z(z)`$ is a rational function, from here one can always factor out $`z_1z_2`$. Thus, the connected component of the polhode in the domain $`S`$ is
$$_{\mathrm{}}=\left\{\frac{Z(z_1)Z(z_2)}{z_1z_2}=0\right\}.$$
(4.11)
Next, for $`\stackrel{~}{\mathrm{}}(U_{})=E_\alpha `$, $`\stackrel{~}{\mathrm{}}^{}(U_{})=0`$, from the addition formula (4.7) we obtain the following simple equation
$$\stackrel{~}{\mathrm{}}^{}(U_1)(E_\alpha \stackrel{~}{\mathrm{}}(U_2))=\stackrel{~}{\mathrm{}}^{}(U_2)(E_\alpha \stackrel{~}{\mathrm{}}(U_1))).$$
(4.12)
Taking squares of both sides and simplifying, we get
$$4(Z_1Z_2)(E_\alpha Z_2)(E_\alpha Z_1)(E_\alpha (Z_1+Z_2)Z_1Z_2+E_\alpha ^2+E_\beta E_\gamma )=0,$$
$$(\alpha ,\beta ,\gamma )=(1,2,3).$$
Then we factor out the term $`(Z_1Z_2)`$ that leads to polhode $`_{\mathrm{}}`$, as well as the product $`(E_\alpha Z_2)(E_\alpha Z_1)`$ that leads to two lines in $`(z_1,z_2)`$-plane and therefore cannot describe the polhode. As a result, we obtain the special generating equation
$$(๐ต(z_1)E_\alpha )(๐ต(z_2)E_\alpha )2E_\alpha ^2E_\beta E_\gamma =0,$$
(4.13)
which defines the special polhode $`_{E_\alpha }`$. For a fixed generic $`z_1`$, this equation has $`N`$ complex solutions for $`z_2`$.
###### Lemma 4.2.
The polhodes $`_{\mathrm{}}`$, $`_{E_\alpha }`$ pass through two vertices of the domain $`S`$.
Proof. Since 6 branch points of $`G`$ are mapped to 4 branch points of $`_2`$, some different finite branch points of $`G`$ are mapped to the same finite branch point on the elliptic curve. Thus, at two vertices of $`S`$, $`๐ต(z_1)=๐ต(z_2)`$ for $`z_1z_2`$ and the polhode $`_{\mathrm{}}`$ passes through these vertices.
Next, at the vertices of $`S`$ one has $`w_1=w_2=0`$, and, in view of the second relation in (4.2), $`\stackrel{~}{\mathrm{}}^{}(U_1)=\stackrel{~}{\mathrm{}}^{}(U_2)=0`$. Hence equation (4.12) is satisfied in all the vertices. On the other hand, at two of the four vertices the condition $`(E_\alpha u_2)(E_\alpha u_1)=0`$ is also satisfied. Since the product $`(E_\alpha u_2)(E_\alpha u_1)`$ was factored out in (4.13), the polhode $`_{E_\alpha }`$ does not pass through the latter two vertices, hence it passes through the other two.
#### Polhodes and Closed Geodesics on an Ellipsoid.
Let $`(\beta ,0)`$ be a finite Weierstrass point on $`G`$. Then under the birational transformation $`(z,w)(\lambda ,\mu )`$ given by (3.5) with $`\alpha =1`$ and $`\beta =b_1`$ (the minimal root of (4.9)) the curve $`G`$ passes to a genus 2 curve
$$\mathrm{\Gamma }=\{\mu ^2=\lambda (\lambda a_1)(\lambda a_2)(\lambda a_3)(\lambda c)\},$$
such that $`a_i`$ and $`c`$ are positive. Thus $`\mathrm{\Gamma }`$ can be regarded as the spectral curve of the geodesic flow on the ellipsoid
$$Q=\left\{\frac{X_1^2}{a_1}+\frac{X_2^2}{a_2}+\frac{X_3^2}{a_3}=1\right\},a_1<a_2<a_3$$
and the corresponding variables
$$\lambda _1=\frac{1}{z_1b_1},\lambda _2=\frac{1}{z_2b_1},$$
(4.14)
are the ellipsoidal coordinates of the moving point on $`Q`$.
In view of Theorem 3.2, under the transformation (3.5) the real asymmetric part of polhode $`_\stackrel{~}{\mathrm{}}_{}`$ describes a closed geodesic on the ellipsoid $`Q`$ in terms of the ellipsoidal coordinates, whereas the whole family of the polhodes gives a one-parametric family of such geodesics that are tangent to one and the same caustic on $`Q`$.
Substituting expressions $`z_1(\lambda _1),z_2(\lambda _2)`$ into the generating equation (4.8), one obtains equation of the geodesic in terms of symmetric functions $`\mathrm{\Sigma }_1=\lambda _1+\lambda _2`$, $`\mathrm{\Sigma }_2=\lambda _1\lambda _2`$ of degree $`2N`$. In view of relations (2.1) for $`n=2`$, the latter can be expressed via the Cartesian coordinates as follows
$`\mathrm{\Sigma }_1`$ $`=a_1+a_3+{\displaystyle \frac{1}{a_1}}(a_2a_1)X_1^2+{\displaystyle \frac{1}{a_3}}(a_2a_3)X_3^2,`$ (4.15)
$`\mathrm{\Sigma }_2`$ $`=a_1a_3+{\displaystyle \frac{a_3}{a_1}}(a_2a_1)X_1^2+{\displaystyle \frac{a_1}{a_3}}(a_2a_3)X_3^2.`$
As a result, one arrives at equation of an algebraic cylinder surface $`๐ฑ_\stackrel{~}{\mathrm{}}_{}`$ of degree 4N in $`^3`$, which cuts out a closed geodesic on $`Q`$. More precisely, one get a family of such surfaces parameterized by $`\stackrel{~}{\mathrm{}}_{}`$.
#### Remark.
Since the equation depends on squares of $`X_i`$ only, such surfaces are symmetric with respect to reflections $`X_iX_i`$. Thus, the complete intersection $`๐ฑ_\stackrel{~}{\mathrm{}}_{}Q`$ consists of a union of closed geodesics that are transformed to each other by these reflections. An example of such intersection is given in Figure 5.4.
As we shall see below, in some cases the equation admit a factorization and the cylinder $`๐ฑ_\stackrel{~}{\mathrm{}}_{}`$ splits in two connected non-symmetric components.
It should be emphasized that the method of polhodes is based on the existence of the second covering $`G_2`$ and the addition law on $`_2`$, so it does not admit a straightforward generalization to a similar description of algebraic closed geodesics on $`n`$-dimensional ellipsoids ($`n>2`$). Indeed, as mentioned in Section 3, in this case $`_2`$ is replaced by an Abelian subvariety $`๐_{n1}`$, for which an algebraic description is not known.
In the sequel we consider in detail polhodes $`_\stackrel{~}{\mathrm{}}_{}`$ and surfaces $`๐ฑ_\stackrel{~}{\mathrm{}}_{}`$ for the 3:1 and 4:1 hyperelliptic tangential covers.
## 5 The 3:1 tangential cover (the Hermite case)
In this case first indicated by Hermite (, see also ), the elliptic curve $`_1`$ in (3.2) is covered by the genus 2 curve
$$G=\left\{w^2=\frac{1}{4}\left(4z^39g_2z27g_3\right)(z^23g_2)\right\}.$$
(5.1)
The latter also covers the second elliptic curve
$$_2=\left\{W^2=\left(4W^3G_2WG_3\right)4(ZE_1)(ZE_2)(ZE_3)\right\},$$
(5.2)
$$G_2=\frac{27}{4}\left(g_2^3+9g_3^2\right),G_3=\frac{243}{8}g_3\left(3g_3^2g_2^3\right)$$
(5.3)
and the covering formulas (4.2) take the form
$$Z=\frac{1}{4}\left(4z^39g_2z9g_3\right),W=\frac{w}{2}\left(4z^23g_2\right).$$
(5.4)
The roots of the polynomial in (5.1) are real iff $`g_2,g_3`$ are real and $`g_2>0`$, $`27g_3^3>g_2^3`$. Then, assuming that $`e_1<e_2<e_3`$, $`E_1<E_2<E_3`$, the following ordering holds
$$b_1=\sqrt{3g_2},\{b_2,b_3,b_4\}=\{3e_1,3e_2,3e_3\},b_5=\sqrt{3g_2},$$
$$E_1=\frac{3}{4}\left(3g_3+\sqrt{3g_2^3}\right),E_2=\frac{3}{4}\left(3g_3\sqrt{3g_2^3}\right),E_3=\frac{9}{2}g_3,$$
(5.5)
and
$$Z(b_1)=E_1,Z(b_5)=E_2,Z(b_2)=Z(b_3)=Z(b_4)=E_3.$$
(5.6)
Now substituting (5.4) into the generating equation (4.8) and taking into account (5.3), we get following family of polhodes
$$M_2(z_1,z_2)\stackrel{~}{\mathrm{}}_{}^2+M_1(z_1,z_2)\stackrel{~}{\mathrm{}}_{}+M_0(z_1,z_2)=0,$$
(5.7)
where
$`M_2`$ $`={\displaystyle \frac{1}{16}}(z_2z_1)^2(9g_24z_1z_24z_1^24z_2^2)^2,`$
$`M_1`$ $`={\displaystyle \frac{729}{16}}g_2^3g_3{\displaystyle \frac{243}{32}}g_2^4\left(z_1+z_2\right)2z_2^3z_1^3\left(z_1^2z_1z_2+z_2^2\right)(z_1+z_2)`$
$`+{\displaystyle \frac{9}{2}}g_2\left(z_1^4+z_2^4z_1z_2^3z_1^3z_2+3z_1^2z_2^2\right)+{\displaystyle \frac{729}{32}}g_3g_2^2\left(4z_1z_2+z_1^2+z_2^2\right)(z_1+z_2)z_1z_2`$
$`{\displaystyle \frac{81}{4}}g_3g_2\left(z_1^4+z_2^4+2z_1z_2^3+2z_1^3z_2\right)+{\displaystyle \frac{27}{32}}g_2^3\left(23z_1z_2+4z_1^2+4z_2^2\right)(z_1+z_2)`$
$`{\displaystyle \frac{81}{8}}g_2^2\left(2z_1^2z_1z_2+2z_2^2\right)(z_1+z_2)z_1z_2+{\displaystyle \frac{9}{2}}g_3\left(z_1^6+z_2^6+4z_1^3z_2^3\right),`$
$`M_0`$ $`={\displaystyle \frac{1}{256}}\left(729g_2^6+4374g_2^5z_1z_2+\mathrm{52\hspace{0.17em}488}g_2^3g_3^2+256z_1^6z_2^6\right)`$
$`+{\displaystyle \frac{\mathrm{10\hspace{0.17em}935}}{128}}g_3g_2^4\left(z_1+z_2\right){\displaystyle \frac{9}{2}}g_3\left(z_1+z_2\right)\left(z_1^2z_1z_2+z_2^2\right)z_1^3z_2^3`$
$`{\displaystyle \frac{9}{2}}g_2\left(z_1^2+z_2^2\right)z_1^4z_2^4+{\displaystyle \frac{6561}{256}}g_3^2g_2^2\left(10z_1z_2+z_1^2+z_2^2\right)`$
$`+{\displaystyle \frac{81}{16}}g_3^2\left(z_1^6+z_2^6+10z_1^3z_2^3\right){\displaystyle \frac{243}{256}}g_2^4\left(8z_1^227z_1z_2+8z_2^2\right)z_1z_2`$
$`+{\displaystyle \frac{81}{16}}g_2^2\left(z_1^4+z_2^4+4z_1^2z_2^2\right)z_1^2z_2^2{\displaystyle \frac{27}{32}}g_2^3\left(27z_1^24z_1z_2+27z_2^2\right)z_1^2z_2^2`$
$`{\displaystyle \frac{729}{32}}g_3^2g_2\left(z_1^4+z_2^4+5z_1z_2^3+5z_1^3z_2\right){\displaystyle \frac{243}{128}}g_3g_2^3\left(20z_1^247z_1z_2+20z_2^2\right)\left(z_1+z_2\right)`$
$`+{\displaystyle \frac{81}{8}}g_3g_2\left(z_1^4+z_2^4z_1z_2^3z_1^3z_2+3z_1^2z_2^2\right)(z_1+z_2)z_1z_2`$
$`{\displaystyle \frac{729}{32}}g_3g_2^2\left(2z_1^2z_1z_2+2z_2^2\right)(z_1+z_2)z_1z_2.`$
#### Reality conditions.
Assume that the parameter $`\stackrel{~}{\mathrm{}}_{}`$ is real and $`(z_1,z_2)`$ range in the square domain $`S`$, namely
$$S=\{3e_1z_13e_2,\mathrm{\hspace{0.33em}3}e_3z_2\sqrt{3g_2}\}.$$
Then, from relations (5.5), (5.6), we conclude that the conditions (4.10) in Lemma 4.1 cannot be satisfied if and only if $`\stackrel{~}{\mathrm{}}_{}`$ varies in the gap $`(\mathrm{};E_1]`$.
#### Limit Polhodes.
When $`\stackrel{~}{\mathrm{}}_{}=\mathrm{}`$, equation (5.7) reduces to
$$9g_24z_1z_24z_1^24z_2^2=0,$$
(5.8)
which can also be obtained directly from (4.11). Thus $`_{\mathrm{}}`$ defines a conic $`๐`$ in $`(z_1,z_2)`$-plane, which passes through 2 vertices of the domain $`S`$.
Next, setting in (4.13) $`\stackrel{~}{\mathrm{}}_{}=E_1=\frac{3}{4}(3g_3\sqrt{3g_2^3})`$ and taking into account expressions (5.3) results in cubic equation
$`16z_1^3z_2^312\sqrt{3g_2^3}\left(z_1^3+z_2^3\right)36g_2z_2z_1\left(z_1^2+z_2^2\right)`$
$`+81g_2^2z_1z_2+27g_2\sqrt{3g_2^3}\left(z_1+z_2\right)+162\sqrt{3g_2^3}g_327g_2^3=0.`$ (5.9)
The corresponding polhode $`_{E_1}`$ also passes through two vertices of $`S`$.
#### Examples of Polhodes in $`S`$.
To illustrate the above polhodes, in (3.2) we choose
$$g_2=3,g_3=0.2.$$
Then the roots of the polynomials $`\frac{1}{4}\left(4z^39g_2z27g_3\right)(z^23g_2)`$ and
$`4Z^3G_2ZG_3`$ are, respectively,
$$(3.0,\mathrm{\hspace{0.33em}2.693},0.201,2.492,3.0),\text{and}(6.3,\mathrm{\hspace{0.33em}0.9},7.2).$$
(5.10)
The domain $`S`$, where the corresponding variables $`w_1,w_2`$ are real, is
$$S=\{2.492z_10.201,2.693z_23.0\}.$$
(5.11)
If the parameter $`\stackrel{~}{\mathrm{}}_{}`$ varies in the interval $`(\mathrm{};7.2)`$, then the real roots of the equation
$$\stackrel{~}{\mathrm{}}_{}=\frac{1}{4}\left(4z^39g_2z9g_3\right)$$
do not fit into the intervals in (5.11). This means that the conditions (4.10) do not hold in domain $`S`$ and the real asymmetric part of the polhode may be non-empty.
Note that if $`\stackrel{~}{\mathrm{}}_{}`$ belongs to other intervals on the real line, some of these conditions are necessarily satisfied, so we exclude this case from consideration.
The graphs of equation (5.7) in the domain (5.11) for two generic values of $`\stackrel{~}{\mathrm{}}_{}`$ are given in Figure 5.1, whereas the graphs of the special polhode $`_{\mathrm{}}`$ given by equation (5.8) and $`_{E_1}`$ given by (5.9) are presented in Figure 5.2.
As seen from Figure 5.1, generic polhodes in $`S`$ intersect generic lines $`z_2=`$const and $`z_1=`$const at 4 and 2 points respectively.
#### Closed geodesics related to 3:1 covering.
Under the the projective transformation (3.5) with $`\beta =b_1=\sqrt{3g_2}`$, the branch points $`\{\sqrt{3g_2},3e_1,3e_2,3e_3,\sqrt{3g_2},\mathrm{}\}`$ of the curve $`G`$ transform to infinity, 4 positive numbers $`\{a_1,a_2,a_3,c\}`$ and zero respectively. Given $`g_2,e_1`$, the parameters $`e_2,e_3`$ of the elliptic curve are defined uniquely:
$$e_2=\frac{e_1}{2}\frac{\sqrt{3}}{6}R,e_3=\frac{e_1}{2}+\frac{\sqrt{3}}{6}R,R=\sqrt{3g_29e_1^2}.$$
Then, assuming that $`a_1<a_2<c<a_3`$, we get
$`a_3`$ $`={\displaystyle \frac{1}{3e_1+B}},a_1={\displaystyle \frac{1}{2B}},`$
$`a_2`$ $`={\displaystyle \frac{1}{3e_3+B}}{\displaystyle \frac{2}{(6e_3B)^2}}\left(2B3e_3\sqrt{3}R\right),`$ (5.12)
$`c`$ $`={\displaystyle \frac{1}{3e_2+B}}{\displaystyle \frac{2}{(6e_3B)^2}}\left(2B3e_3+\sqrt{3}R\right),`$
where $`B=\sqrt{3g_2}`$. As a result, the four parameters $`a_1<a_2<c<a_3`$ are uniquely defined by $`g_2`$ (or $`B`$) and $`e_1`$.
Now we apply the transformation (4.14) with $`b_1=\sqrt{3g_2}`$ to the polhode (5.7). This yields an equation of a closed geodesic on $`Q`$ written in terms of the symmetric functions $`\mathrm{\Sigma }_1=\lambda _1+\lambda _2`$, $`\mathrm{\Sigma }_2=\lambda _1\lambda _2`$ of the ellipsoidal coordinates. (In fact, one obtains a family of such geodesics parameterized by $`\stackrel{~}{\mathrm{}}_{}`$.) Then, making the substitution (4.15) one obtains the equation of the cylinder cutting surface $`๐ฑ_\stackrel{~}{\mathrm{}}_{}`$ in terms of squares of the Cartesian coordinates $`X_1,X_3`$. For a generic parameter $`\stackrel{~}{\mathrm{}}_{}`$ this equation has degree 12, it is quite tedious and we do not give it here. However, the structure of a generic polhode in $`S(z_1,z_2)`$ and the correspondence between the sets $`\{3e_1,3e_2,3e_3,\sqrt{3g_2}\}`$ and $`\{a_1,a_2,c,a_3\}`$ is already sufficient to give a complete qualitative description of the geodesic on $`Q`$.
Namely, let $`_c`$ be a ring on $`Q`$ bounded by the two connected components of the caustic $`QQ_c`$ and $`\rho =k:l`$ be the quotient of the numbers of complete rotations performed by a closed geodesics in lateral and meridional directions on the ring respectively (the rotation number).
###### Theorem 5.1.
1). Under the assumption $`a_1<a_2<c<a_3`$, the geodesic corresponding to a generic polhode (5.7) or to the special polhodes is located in the ring $`_c`$ between planes $`X_3=\pm h`$, $`h<\sqrt{a_3}`$ and has rotation number $`\rho =2:1`$. It touches the caustic $`QQ_c`$ at 2 points and has one self-intersection.
2). Under the assumption $`a_1<c<a_2<a_3`$, the geodesic is located in the ring $`_c`$ between planes $`X_1=\pm h`$, $`h<\sqrt{a_1}`$ and has rotation number $`1:2`$. It touches the caustic $`QQ_c`$ at 4 points and has no self-intersections.
In both cases the geodesic is either a 2-fold covering of the real asymmetric part of a generic polhode $`_\stackrel{~}{\mathrm{}}_{}`$ or a 4-fold covering of that of the special polhodes.
Note that the self-intersection point of the polhode does not correspond to the self-intersection point of the corresponding closed geodesic.
Sketch of Proof of Theorem 5.1. Under the projective transformation (4.14), a polhode $`_\stackrel{~}{\mathrm{}}_{}(z_1,z_2)`$ is mapped to a polhode $`\stackrel{~}{}_\stackrel{~}{\mathrm{}}_{}`$ in $`^2=(\lambda _1,\lambda _2)`$, which is tangent to lines $`\lambda _j=a_i`$, $`i=1,2,3`$ and $`\lambda _j=c`$. In view of relations (2.1), the point of tangency of $`\stackrel{~}{}_\stackrel{~}{\mathrm{}}_{}`$ to the line $`\lambda _j=a_i`$ corresponds to the moment when the geodesic $`X(s)`$ on $`Q`$ crosses the plane $`X_i=0`$, and the tangency to the line $`\lambda _j=c`$ corresponds to the tangency of $`X(s)`$ to the caustic $`QQ_c`$. Estimating ordering and number of the tangencies of $`\stackrel{~}{}_\stackrel{~}{\mathrm{}}_{}`$ in the cases $`a_2<c`$ and $`c<a_2`$, one arrives at the statements of the theorem. $`\overline{)}`$
#### An Example of a Generic Closed Geodesic.
For the above numerical choice of $`g_2,g_3`$ one gets $`B=3`$, $`e_1=0.83054`$, $`e_2=0.067069`$, $`e_3=0.89761`$ and the formulas (5.12) (or the images of the values in (5.10)) yield
$$a_1=0.16667,a_2=0.1776,c=0.3579,a_3=1.96703.$$
(This means that the corresponding ellipsoid is almost โprolateโ<sup>7</sup><sup>7</sup>7 In all our numeric examples some of the branch points of the curve $`\mathrm{\Gamma }`$ are rather close to each other. Apparently, this phenomenon is unavoidable and is due to the projective transformation (4.14)..) Projections of the intersection $`Q๐ฑ_{\mathrm{}}`$ onto $`(X_1,X_3)`$\- and $`(X_1,X_2)`$-planes for $`\stackrel{~}{\mathrm{}}_{}=11`$ are given in Figure 5.3.
One can see that this intersection actually consists of four closed geodesics obtained from each other by reflections $`(X_1,X_2)(\pm X_1,\pm X_2)`$. Each geodesics has the only self-intersection point at $`X_3=0`$ and corresponds to the polhode in Figure 5.1 (b) which is passed two times.
It is natural to conjecture that the four geodesics are real parts of one and the same spatial elliptic curve which are obtained from each other via translations by elements of a finite order subgroup of the curve.
#### Special Geodesic for $`_{\mathrm{}}`$.
In the special case $`\stackrel{~}{\mathrm{}}_{}=\mathrm{}`$ the equation of the surface $`๐ฑ_{\mathrm{}}`$ simplifies drastically and admits the following factorization
$$\left(\alpha (X_1\gamma )^2+\beta X_3^2\delta \right)\left(\alpha \left(X_1+\gamma \right)^2+\beta X_3^2\delta \right)=0,$$
(5.13)
where
$`\alpha `$ $`={\displaystyle \frac{\sqrt{2}}{9}}\left(B6e_1\right)(\left(2B+3e_1\right)\sqrt{3}R),`$ (5.14)
$`\beta `$ $`={\displaystyle \frac{1}{3\sqrt{6}}}\left(B+3e_1\right)\left(R3\sqrt{3}e_1\right),`$
$`\gamma `$ $`={\displaystyle \frac{2}{9\alpha }}\left(R+3\sqrt{3}e_1\right)R\sqrt{B},`$
$`\delta `$ $`={\displaystyle \frac{B^2}{3\sqrt{2}}}{\displaystyle \frac{\left(R+3\sqrt{3}e_1\right)^2}{\left(\sqrt{3}R\left(3e_1+2B\right)\right)\left(B6e_1\right)}}.`$
and, as above, $`B=\sqrt{3g_2}`$, $`R=\sqrt{3g_29e_1^2}`$.
Equation (5.13) defines a union of two elliptic cylinders in $`^3`$ that are transformed to each other by mirror symmetry with respect to the plane $`X_1=0`$. It appears that each cylinder is tangent to the ellipsoid $`Q`$ at a point $`(X_2=X_3=0)`$ and cuts out a closed geodesic with the only self-intersection at this point. As a result, the special closed geodesic on $`Q`$ related to the polhode (5.8) is defined by its intersection with just a quadratic surface defined by one of the two factors in (5.13)).
#### Remark.
Note that due to the self-intersection, the special geodesic in $`^3`$ ($`^3`$) is a rational algebraic curve and not an elliptic one, as the intersection of two generic quadrics. It admits parameterization
$$X_1=d+h_1\mathrm{cos}(2\nu ),X_2=h_2\mathrm{sin}(2\nu ),X_3=h_3\mathrm{sin}\nu ,\nu ,$$
$`h_i,d`$ being certain constants<sup>8</sup><sup>8</sup>8Here the parameter $`\nu `$ is not a linear function of time $`t`$ or the rescaling parameter $`s`$ in (2.2)..
On the other hand, in the phase space $`(X,\dot{X})`$ the corresponding periodic solution has no self-intersections and represents an elliptic curve. Indeed, in view of formulas (2.1), the latter can be regarded as a 4-fold covering of the rational special polhode (5.8). The covering has simple ramifications at 8 points that are projected to two vertices of the domain $`S`$ and two vertices of the symmetric domain $`S^{}`$ obtained by reflection with respect to the diagonal $`z_1=z_2`$. Then, according to the RiemannโHurwitz formula (see, e.g., ), the covering has genus one. The projection $`^6=\{(X,\dot{X})\}^3=\{X\}`$ maps two different points of the elliptic solution to the self-intersection point on $`Q`$.
For the above values of the parameters $`a_1,a_2,a_3,c`$ the 3D graph of the special geodesic is shown in Figure 5.4.
#### Special Geodesic for $`_{E_1}`$.
Applying the transformation (4.14) with $`b_1=\sqrt{3g_2}`$ to the special polhode (5.9) and making the substitution (4.15) we arrive at a sextic surface in $`^3`$ given by equation
$`256B^3f_1^3X_1^6\left(48B^3f_1f_2^2+432Be_1^2f_1f_2^2+288B^2e_1f_1f_2^2\right)X_1^2X_3^4`$
$`+\left(2304B^3e_1f_1^2192B^4f_1^26912B^2e_1^2f_1^2\right)X_1^4`$
$`+\left(192B^4f_1f_2+5184Be_1^3f_1f_2864B^3e_1f_1f_22592B^2e_1^2f_1f_2\right)X_1^2X_3^2`$
$`+\left(4B^3f_2^3108e_1^3f_2^3108Be_1^2f_2^336B^2e_1f_2^3\right)X_3^6`$
$`+\left(3B^5f_2\mathrm{11\hspace{0.17em}664}e_1^5f_227B^4e_1f_2+1296B^2e_1^3f_2108B^3e_1^2f_2\right)X_3^2`$
$`+\left(864B^4e_1f_1\mathrm{46\hspace{0.17em}656}Be_1^4f_136B^5f_1+\mathrm{31\hspace{0.17em}104}B^2e_1^3f_17776B^3e_1^2f_1\right)X_1^2`$
$`+\left(1944e_1^4f_2^212B^4f_2^2+648Be_1^3f_2^2144B^3e_1f_2^2324B^2e_1^2f_2^2\right)X_3^4`$
$`\left(192B^3f_1^2f_2+576B^2e_1f_1^2f_2\right)X_1^4X_3^2=0,`$ (5.15)
where
$$f_1=\left(9e_1^2\frac{7}{4}B^2+BR\sqrt{3}\right),f_2=\left(27e_1^2\frac{3}{2}B^29Be_1+BR\sqrt{3}+3Re_1\sqrt{3}\right).$$
It cuts out a pair of closed geodesic on $`Q`$ that are transformed to each other by mirror symmetry with respect to the plane $`X_2=0`$. Both geodesics have a 3D shape similar to that in Figure 5.4, each of them has the only self-intersection point for $`(X_1=X_3=0)`$.
Note, however, that in contrast to quartic equation (5.13), the sextic polynomial in (5.15) does not admit a factorization, hence none of the above geodesics can be represented as the intersection of $`Q`$ with a quadratic or a cubic cylinder.
For the above choice of moduli $`g_2,g_3`$ and the parameters $`a_i,c`$ the projection of the sextic surface and the corresponding geodesics onto $`(X_1,X_3)`$-plane are given in Figure 5.5.
#### Remark.
As follows from the above considerations, under the condition $`a_1<a_2<c<a_3`$ all the real closed geodesics of the one-parametric family have one self-intersection point on the equator $`\{X_3=0\}Q`$, and as the parameter $`\stackrel{~}{\mathrm{}}_{}`$ ranges from $`\mathrm{}`$ to $`E_1`$, this point varies from the $`X_1`$-axis to $`X_2`$-axis.
The case $`a_1<c<a_2<a_3`$ will be illustrated in detail elsewhere.
## 6 The case of 4:1 tangential covering
This case was originally studied by Darboux and later appeared in paper in connection with new elliptic solutions of the KdV equation (see also ). Namely, the genus 2 curve $`G`$
$$w^2=(z6e_1)\underset{l=1}{\overset{4}{}}(zz_l),$$
(6.1)
with
$`z_{1,2}`$ $`=e_3+2e_2\pm \sqrt{28e_2^2+76e_2e_3+40e_3^2},`$
$`z_{3,4}`$ $`=e_2+2e_3\pm \sqrt{28e_3^2+76e_2e_3+40e_2^2},`$
is a 4-fold cover of the curve (3.2). It also covers second elliptic curve
$`_2=\left\{W^2=4(ZE_1)(ZE_2)(ZE_3)\right\}`$, such that
$`E_1`$ $`=3\left(e_2e_3\right)^36\left(e_2e_1\right)(5e_12e_3)^2,`$
$`E_2`$ $`=6\left(e_2e_3\right)^3+3\left(e_2e_1\right)(5e_12e_3)^2,`$ (6.2)
$`E_3`$ $`=3\left(e_2e_3\right)^3+3\left(e_2e_1\right)(5e_12e_3)^2`$
as described by one of the formulas
$`Z`$ $`=E_1+F_\alpha (z),F_\alpha ={\displaystyle \frac{9}{4}}{\displaystyle \frac{\left(z^23e_\alpha z24\left(e_\beta ^2+e_\gamma ^2\right)51e_\beta e_\gamma \right)^2}{z6e_\alpha }},`$ (6.3)
$`W`$ $`=w{\displaystyle \frac{dZ(z)}{dz}},`$
where $`(\alpha ,\beta ,\gamma )`$ is a circular permutation of $`(1,2,3)`$. In the sequel we assume
$`\alpha =1,\beta =2,\gamma =3`$.
Substituting projection formulas (6.3) to the generating equation (4.8) one obtains equation of generic polhodes of degree 16, which is much more tedious than the family (5.7) for the 3:1 cover, so we do not give it here.
#### The special polhodes for $`\stackrel{~}{\mathrm{}}_{}=\mathrm{}`$ and $`\stackrel{~}{\mathrm{}}_{}=E_1`$.
Substituting projection formulas (6.3) to the special generating equation (4.11), we obtain algebraic equation of degree 4,
$`6e_1\left(z_1^3+z_2^3\right)36e_1^2\left(z_1^2+z_2^2\right)z_1^2z_2^2+\left(12e_1(z_1+z_2)z_1^2z_2^2\right)z_2z_1`$
$`+(6e_2e_3+3e_1^2)z_1z_2+\left(54e_1^3612e_1e_2e_3288e_1e_2^2288e_1e_3^2\right)(z_1+z_2)`$
$`+9\left(17e_2e_3+8e_2^2+8e_3^2\right)\left(17e_2e_3+12e_1^2+8e_2^2+8e_3^2\right)=0.`$ (6.4)
Next, substituting (6.3) into the special generating equation (4.13) with $`E_\alpha =E_1`$ and taking into account relation $`E_1+E_2+E_3=0`$, one gets the following equation of degree 8
$$81z_1^4z_2^4486e_1(z_2^3z_1^4+z_2^4z_1^3)(3159e_2^2+6804e_2e_3+3159e_{3}^{}{}_{}{}^{2})(z_2^2z_1^4+z_2^4z_1^2)$$
$$+2916e_1^2z_1^3z_2^3+2916e_1\xi _1(z_2z_1^4+z_2^4z_1)+1458e_1\xi _2(z_2^2z_1^3+z_2^3z_1^2)$$
$$+729\xi _1^2(z_1^4+z_2^4)8748\xi _1e_1^2(z_2z_1^3+z_2^3z_1)+729\xi _2^2z_1^2z_2^24374e_1\xi _1^2(z_1^3+z_2^3)$$
$$4374e_1\xi _1\xi _2(z_2z_1^2+z_2^2z_1)2187\xi _2\xi _1^2(z_1^2+z_2^2)+16(104976e_{2}^{}{}_{}{}^{6}+656100e_{2}^{}{}_{}{}^{5}e_3$$
$$+\frac{6725025}{4}e_{2}^{}{}_{}{}^{4}e_{3}^{}{}_{}{}^{2}+\frac{4520529}{2}e_{2}^{}{}_{}{}^{3}e_{3}^{}{}_{}{}^{3}+\frac{6725025}{4}e_{2}^{}{}_{}{}^{2}e_{3}^{}{}_{}{}^{4}+656100e_2e_{3}^{}{}_{}{}^{5}+104976e_{3}^{}{}_{}{}^{6}2E_{2}^{}{}_{}{}^{2}$$
$$5E_2E_32E_{3}^{}{}_{}{}^{2})z_1z_2+\frac{3}{8}e_1(1119744e_{2}^{}{}_{}{}^{6}+7138368e_{2}^{}{}_{}{}^{5}e_3+18528264e_{2}^{}{}_{}{}^{4}e_{3}^{}{}_{}{}^{2}$$
$$+25021467e_{2}^{}{}_{}{}^{3}e_{3}^{}{}_{}{}^{3}+18528264e_{2}^{}{}_{}{}^{2}e_{3}^{}{}_{}{}^{4}+7138368e_2e_{3}^{}{}_{}{}^{5}+1119744e_{3}^{}{}_{}{}^{6}+32E_{2}^{}{}_{}{}^{2}$$
$$+80E_2E_3+32E_{3}^{}{}_{}{}^{2})(z_1+z_2)+52225560e_2^6e_3^2+107298594e_2^5e_3^3+\frac{2165451489}{16}e_2^4e_3^4$$
$$+14276736e_{2}^{}{}_{}{}^{7}e_3144e_2e_3E_{3}^{}{}_{}{}^{2}72e_{3}^{}{}_{}{}^{2}E_{2}^{}{}_{}{}^{2}180e_{3}^{}{}_{}{}^{2}E_2E_372e_{3}^{}{}_{}{}^{2}E_{3}^{}{}_{}{}^{2}$$
$$+107298594e_{2}^{}{}_{}{}^{3}e_{3}^{}{}_{}{}^{5}+52225560e_{2}^{}{}_{}{}^{2}e_{3}^{}{}_{}{}^{6}+14276736e_2e_3^7+1679616e_{3}^{}{}_{}{}^{8}72e_{2}^{}{}_{}{}^{2}E_{2}^{}{}_{}{}^{2}$$
$$180e_{2}^{}{}_{}{}^{2}E_2E_372e_{2}^{}{}_{}{}^{2}E_{3}^{}{}_{}{}^{2}144e_2e_3E_{2}^{}{}_{}{}^{2}360e_2e_3E_2E_3+1679616e_{2}^{}{}_{}{}^{8}=0,$$
(6.5)
$$\xi _1=8e_{2}^{}{}_{}{}^{2}+17e_2e_3+8e_{3}^{}{}_{}{}^{2},\xi _2=13e_{2}^{}{}_{}{}^{2}+28e_2e_3+13e_{3}^{}{}_{}{}^{2},$$
#### Examples of polhodes in $`S`$.
To illustrate the above polhodes, in the first elliptic curve (3.2) we choose
$$e_1=2,e_2=1,e_3=3,\text{ so that }E_1=1728,E_2=1152,E_3=576,$$
and the roots of the polynomial in (6.1) become
$$(12.0,11.649,3.0,\mathrm{\hspace{0.33em}13.0},\mathrm{\hspace{0.33em}13.649}).$$
(6.6)
As a result, the square domain $`S(z_1,z_2)`$ where the corresponding variables $`w_1,w_2`$ are real is
$$S=\{11.649z_13.0,13.0z_213.649\}.$$
(6.7)
If the parameter $`\stackrel{~}{\mathrm{}}_{}`$ ranges in the interval $`(\mathrm{};E_1=1728)`$, then the conditions (4.10) do not hold in $`S`$, hence the real asymmetric part of the polhode is non-empty.
Graphs of polhodes in the domain (6.7) for a generic value of $`\stackrel{~}{\mathrm{}}_{}(\mathrm{};1728)`$ is given in Figure 6.2, whereas the special polhodes defined by (6.4) and (6.5) are shown in Figure 6.2.
As seen from Figure 6.2, the generic polhodes in $`S`$ intersect generic lines $`z_2=`$const and $`z_1=`$const at 6 and 2 points respectively.
#### Special Closed Geodesic for $`_{\mathrm{}}`$.
Assuming, as above, $`e_1<e_2<e_3`$, we conclude that $`b_1=6e_1`$ is the minimal root of the hyperelliptic polynomial in (6.1). Setting this value into the transformation (4.14), we get
$$z_1=\frac{6e_1\lambda _1+1}{\lambda _1},z_2=\frac{6e_1\lambda _2+1}{\lambda _2}.$$
The latter substitution transforms equation (6.4) of polhode $`_{\mathrm{}}`$ to
$`\left(102e_2e_3117e_1^2+48e_2^2+48e_3^2\right)\lambda _1^2\lambda _2^2\lambda _1^2\lambda _2^2\lambda _1\lambda _2`$
$`18e_1\left(\lambda _1+\lambda _2\right)\lambda _2\lambda _1+9\left(17e_2e_36e_1^2+8e_2^2+8e_3^2\right)^2\lambda _1^3\lambda _2^3=0,`$ (6.8)
which describes the closed geodesic in ellipsoidal coordinates on $`Q`$. Next, assuming that $`a_1<a_2<c<a_3`$, we find
$$a_1=\frac{1}{6e_1+e_2+2e_3+\sqrt{28e_3^2+76e_2e_3+40e_2^2}},a_2=\frac{1}{6e_1+e_2+2e_3\sqrt{28e_3^2+76e_2e_3+40e_2^2}},$$
$$a_3=\frac{1}{6e_1+e_3+2e_2\sqrt{28e_2^2+76e_2e_3+40e_3^2}},c=\frac{1}{6e_1+e_3+2e_2+\sqrt{28e_2^2+76e_2e_3+40e_3^2}}.$$
Applying formulas (4.15) and simplifying, one obtains equation of cylinder surface $`๐ฑ_{\mathrm{}}`$ of degree 6 in coordinates $`(X_1,X_3)`$, which admits the factorization
$$F_{}(X_1,X_3)F_+(X_1,X_3)=0,$$
$$F_\pm =\pm h_{30}X_1^3+h_{03}X_3^3+h_{21}X_1^2X_3\pm h_{12}X_1X_3^2\pm h_{10}X_1+h_{01}X_3,$$
(6.9)
where $`h_{ij}`$ are rather complicated expressions of $`e_2,e_3`$, so we do not give them here.
Thus $`๐ฑ_{\mathrm{}}`$ is a union of two cubic cylinders $`๐ฑ_{\mathrm{}}`$, $`๐ฑ_\mathrm{}+`$ that are obtained from each other by mirror symmetry with respect to the plane $`X_1=0`$.
#### Example of a closed geodesic for $`_{\mathrm{}}`$.
Under the above assumption $`a_1<a_2<c<a_3`$, the values (6.6) lead to numbers
$$a_1=3.69877\times 10^2,a_2=0.04,c=0.111111,a_3=2.84981$$
and the equation of the cylinder $`๐ฑ_\mathrm{}+`$ becomes
$$0.301917X_1^30.470556X_1X_3^2+0.652885X_1^2X_3+1.87228X_10.303831X_3+0.113051X_3^3=0.$$
(6.10)
Its projection onto $`(X_1,X_3)`$-plane and the 3D graph of the corresponding geodesic are shown in Figure 6.3 <sup>9</sup><sup>9</sup>9The graph was actually produced by direct numeric integration of the corresponding geodesic equation with initial conditions prescribed by (6.10). The cylinder is tangent to the ellipsoid $`Q`$ at 2 points with $`X_2=0`$, which implies that the geodesic has two self-intersection points.
#### Remark.
The special geodesic can be regarded as a two-fold covering of the plane algebraic curve $`_\pm =\{F_\pm (X_1,X_3)=0\}`$, which is ramified at two points of transversal intersection of $`_\pm `$ with the ellipse $`\{X_1^2/a_1+X_3^2/a_3=1\}`$. (There is no ramification at the two points of tangency with the ellipse.) Using explicit expressions of (6.9), one can show that for any value of $`e_2,e_3`$ the genus of $`_\pm `$ equals zero. Hence, according to the RiemannโHurvitz formula, the special closed geodesic is a rational curve.
However, similarly to the special geodesic for the 3:1 covering, in the phase space $`(X,\dot{X})`$ the periodic solution corresponding to $`_{\mathrm{}}`$ has no self-intersections and represents a connected component of real part of an elliptic curve.
#### General closed geodesics.
For a generic parameter $`\stackrel{~}{\mathrm{}}_{}`$ the equation of the cutting cylinder $`๐ฑ_\stackrel{~}{\mathrm{}}_{}`$ has degree 16 and its projection onto $`(X_1,X_3)`$-plane looks quite tangled: it includes 2 intersecting closed geodesics obtained from each other by reflections $`X_1X_1`$. The structure of the real asymmetric part of generic and the special polhodes implies the following behavior under the condition $`a_1<a_2<c<a_3`$: all the closed geodesics of the family have two centrally symmetric self-intersection points and as the parameter $`\stackrel{~}{\mathrm{}}_{}`$ ranges from $`\mathrm{}`$ to $`E_1`$, these points vary from the plane $`\{X_2=0\}`$ to $`\{X_1=0\}`$. All the geodesics have rotation number $`\rho =3:1`$.
Similarly to the case of the 3-fold tangential covering, one can consider the ordering
$`a_1<c<a_2<a_3`$, which leads to generic and special closed geodesics on $`Q`$ without self-intersections.
### Conclusion
In this paper we proposed a simple method of explicit constructing families of algebraic closed geodesics on triaxial ellipsoids, which is based on properties of tangential coverings of an elliptic curve and the addition theorem for elliptic functions. We applied the method to the cases of 3- and 4-fold coverings and gave concrete examples of algebraic surfaces that cut such closed geodesics. The latter coincide with those obtained by direct numeric integration of the geodesic equation. This serves as an ultimate proof of correctness of the method. Depending on how one chooses the caustic parameter $`c`$ in the interval $`(a_1,a_3)`$, the closed geodesics may or may not have self-intersections.
Thus, our approach can be regarded as a useful application of the WeierstrassโPoincarรฉ reduction theory. One must only know explicit covering formulas (4.2), as well as expressions for $`z`$-coordinates of the finite branch points of the genus 2 curve $`G`$ in terms of the moduli $`g_2,g_3`$. To our knowledge, until now such expressions are calculated only for $`N8`$.
Since the method essentially uses the algebraic addition law on the second elliptic curve $`_2`$, it does not admit a straightforward generalization to similar description of algebraic closed geodesics on $`n`$-dimensional ellipsoids: as mentioned in Section 3, in this case $`_2`$ is replaced by an Abelian subvariety $`๐_{n1}`$, for which an algebraic description is not known.
On the other hand, one should not exclude the existence of algebraic closed geodesics on ellipsoids related to other type of doubly periodic solutions of the KdV equation, e.g., elliptic not $`x`$, but in $`t`$-variable. This is expected to be a subject of a separate study.
Our approach can equally be applied to describe elliptic solutions of other integrable systems linearized on two-dimensional hyperelliptic Jacobians or their coverings.
### Acknowledgments
I am grateful to A. Bolsinov, L. Gavrilov, V. Enolskii, E. Previato, and A. Treibich for stimulating discussions, as well as to A. Perelomov for some important suggestions during preparation of the manuscript and indicating me the reference . I also thank the referees for their valuable remarks that helped to improve the text.
The support of grant BFM 2003-09504-C02-02 of Spanish Ministry of Science and Technology is gratefully acknowledged.
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# Noncommutative General Relativity
## 1 Introduction
General Relativity is a very successful theory when it comes to describe macroscopic effects of gravitation. However, it is widely believed that an unification of Quantum Mechanics and General Relativity requires a short distance modification of spacetime. It can be shown that classical General Relativity considered together with Quantum Mechanics implies the existence of a fundamental length . A class of models that incorporate the notion of a fundamental length in gauge theories are gauge theories formulated on noncommutative spaces.
It is a challenge to formulate General Relativity on noncommutative spaces and there are thus different approaches in the literature. In for example a deformation of Einsteinโs gravity was studied using a construction based on gauging the noncommutative SO(4,1) de Sitter group and the Seiberg-Witten map with subsequent contraction to ISO(3,1). Most recently another construction of a noncommutative gravitational theory was proposed based on a twisted Poincarรฉ algebra . These approaches although mathematically consistent, are not minimal formulations of Einsteinโs General Relativity on noncommutative spaces. The main problem in formulating a theory of gravity on noncommutative manifolds is that it is difficult to implement symmetries such as general coordinate covariance and local Lorentz invariance and to define derivatives which are torsion-free and satisfy the metricity condition.
Similar obstacles appear in constructing models of particle physics on flat spacetime with canonical noncommutativity defined by the algebra $`[\widehat{x}^a,\widehat{x}^b]=i\theta ^{ab}`$ ($`\theta ^{ab}`$ is constant and antisymmetric). Indeed, it turns out to be rather difficult to implement most symmetries particle physicists are so familiar with. In particular, Lorentz invariance is explicitly violated by the noncommutative algebra. However, it has been shown in that there is another exact symmetry, the noncommutative Lorentz invariance, based on the usual Lorentz algebra so(3,1) which is undeformed. Another example is the implementation of noncommutative local gauge theories. A formulation of noncommutative gauge theories within the enveloping algebra approach has been proposed in . The fields taken in the enveloping algebra are expanded in term of a series in $`\theta `$. Each of the terms of this series is however a function of the classical variables . The number of degrees of freedom is thus finite and the same as in the corresponding commutative gauge theory.
In this work, based partially on the above achievements in implementing symmetries on flat noncommutative spacetimes, we would like to propose a theory of General Relativity on curved spacetimes with canonical noncommutativity. We shall use the tetrad approach to General Relativity. This formalism applied to noncommutative General Relativity allows to follow closely the usual construction of noncommutative gauge theories. This requires to implement two gauge symmetries: local Lorentz transformations which can be seen as a local gauge theory based on the algebra so(3,1) for the spin connection field and general coordinate transformations which are inhomogenous translations with the tetrad as a gauge field.
The gauging of noncommutative so(3,1) algebra is only possible if the corresponding gauge field, the spin connection, is assumed to be in the enveloping algebra. Hence, to implement local Lorentz invariance we follow the approach developed in . The invariance under the general coordinate transformations, however, is explicitly violated by the canonical noncommutative algebra. Nevertheless, we find a restricted class of coordinate transformations which preserve the canonical structure. It turns out that this transformations correspond to volume-preserving diffeomorphisms. Thus, the basic new ideas for constructing a theory of noncommutative General Relativity are to formulate local Lorentz invariance by gauging so(3,1) within the enveloping algebra approach and to introduce volume-preserving coordinate transformations in place of general coordinate transformations, which is indeed an exact symmetry of the canonical noncommutative spacetime.
## 2 Noncommutative General Relativity
Let us start from a noncommutative spacetime and assume that the coordinates fulfill canonical commutation relations:
$`[\widehat{x}^\mu ,\widehat{x}^\nu ]=i\theta ^{\mu \nu }.`$ (1)
Obviously, the commutator (1) explicitly violates general coordinate covariance since $`\theta ^{\mu \nu }`$ is constant in all reference frames. However, we can identify a subclass of general coordinate transformations,
$$\widehat{x}^\mu =\widehat{x}^\mu +\widehat{\xi }^\mu (\widehat{x}),$$
(2)
which are compatible with the algebra given by (1). The hat on the function $`\widehat{\xi }(\widehat{x})`$ indicates that it is in the enveloping algebra. Under the change of coordinates (2) the commutator (1) transforms as:
$`[\widehat{x}^\mu ,\widehat{x}^\nu ]`$ $`=`$ $`\widehat{x}^\mu \widehat{x}^\nu \widehat{x}^\nu \widehat{x}^\mu =i\theta ^{\mu \nu }+[\widehat{x}^\mu ,\widehat{\xi }^\nu ]+[\widehat{\xi }^\mu ,\widehat{x}^\nu ]+๐ช(\widehat{\xi }^2)`$ (3)
Requiring that $`\theta `$ remains constant yields the following partial differential equations:
$`\theta ^{\mu \alpha }\widehat{}_\alpha \widehat{\xi }^\nu (\widehat{x})=\theta ^{\nu \beta }\widehat{}_\beta \widehat{\xi }^\mu (\widehat{x}).`$ (4)
A nontrivial solution to this condition can be easily found:
$$\widehat{\xi }^\mu (\widehat{x})=\theta ^{\mu \nu }\widehat{}_\nu \widehat{f}(\widehat{x}),$$
(5)
where $`\widehat{f}(\widehat{x})`$ is an arbitrary field. This noncommutative general coordinate transformation corresponds to the following classical transformation: $`\widehat{\xi }^\mu (x)=\theta ^{\mu \nu }_\nu \widehat{f}(x)`$. The Jacobian of this restricted coordinate transformations is equal to 1, meaning that the volume element is invariant: $`d^4x^{}=d^4x`$. The version of General Relativity based on volume-preserving diffeomorphism is known as the unimodular theory of gravitation . Thus we came to the conclusion that symmetries of canonical noncommutative spacetime naturally lead to the noncommutative version of unimodular gravity.
Now we need to implement two gauge symmetries mentioned above. A noncommutative gauge transformation $`\widehat{\mathrm{\Lambda }}(\widehat{x})`$ valued in the iso(3,1) Lie algebra can be decomposed using the generators of the inhomogeneous translations $`p_\mu =i_\mu `$, which are anti-Hermitian, and the generators of the Local Lorentz algebra so(3,1) $`\mathrm{\Sigma }_{ab}`$, which are Hermitian. One finds
$`\widehat{\mathrm{\Lambda }}(\widehat{x})=\widehat{\xi }(\widehat{x})+\widehat{\mathrm{\Lambda }}(\widehat{x})=\widehat{\xi }^\mu (\widehat{x})p_\mu +{\displaystyle \frac{1}{2}}\widehat{\lambda }^{ab}(\widehat{x})\mathrm{\Sigma }_{ab},`$ (6)
where $`\widehat{\xi }^\mu `$ is subject to the constraint (5). Note that $`p_\mu `$ acts on the coordinates and functions, including $`\widehat{\lambda }^{ab}`$, and $`a,b,\mathrm{}`$ run over the tangent space indices. As in , the algebra of generators is undeformed. It is easy to verify that the commutator of two noncommutative gauge transformations $`[\widehat{\mathrm{\Lambda }}_1(\widehat{x}),\widehat{\mathrm{\Lambda }}_2(\widehat{x})]`$ is in general not a noncommutative gauge transformation if the transformations are Lie algebra valued. As in the Yang-Mills case, the solution is to assume that the noncommutative gauge transformations are in the enveloping algebra. Let us introduce a noncommutative vector potential which corresponds to the noncommutative gauge transformation (6)
$`\widehat{A}_a(\widehat{x})=(\widehat{D}_a)=i\widehat{E}_a^\mu (\widehat{x})p_\mu +{\displaystyle \frac{i}{2}}\widehat{\omega }(\widehat{x})_a^{bc}\mathrm{\Sigma }_{bc}`$ (7)
where $`\widehat{E}_a^\mu (\widehat{x})`$ are the components of the noncommutative tetrad $`\widehat{E}_a(\widehat{x})`$ , i.e. the gauge fields corresponding to general coordinate transformations and $`\widehat{\omega }(\widehat{x})_a^{bc}`$ are the spin connections fields associated with local Lorentz invariance. Note that $`\widehat{A}_a(\widehat{x})`$ plays a dual role. It can be viewed as a covariant derivative as well. It is also worth noticing that $`\widehat{E}_a=\widehat{E}_a^\mu \widehat{}_\mu =\widehat{}_a`$, which implies that the noncommutative tetrad is mapped trivially on the commutative one: $`\widehat{E}_a=e_a`$ to all orders in $`\theta `$.
Let us now assume that the gauge transformations and the spin connection field are in the enveloping algebra:
$`\widehat{\mathrm{\Lambda }}=\mathrm{\Lambda }(x)+\mathrm{\Lambda }^{(1)}(x,\omega _a)+๐ช(\theta ^2),`$ (8)
and
$`\widehat{\omega }_a=\omega _a(x)+\omega _a^{(1)}(x,\omega _a)+๐ช(\theta ^2),`$ (9)
respectively, with $`\mathrm{\Lambda }(x)=\xi ^\mu (x)p_\mu +\frac{1}{2}\lambda ^{ab}(x)\mathrm{\Sigma }_{ab}`$ and $`\omega _a(x)=\frac{1}{2}\omega _a^{bc}\mathrm{\Sigma }_{bc}`$. We require that the commutator of two noncommutative gauge transformations with $`\widehat{\mathrm{\Lambda }}_1`$ and $`\widehat{\mathrm{\Lambda }}_2`$ be a gauge transformation $`\widehat{\mathrm{\Lambda }}_{\widehat{\mathrm{\Lambda }_1\times \mathrm{\Lambda }_2}}`$:
$`\left(\widehat{\delta }_{\widehat{\mathrm{\Lambda }}_1}\widehat{\delta }_{\widehat{\mathrm{\Lambda }}_2}\widehat{\delta }_{\widehat{\mathrm{\Lambda }}_2}\widehat{\delta }_{\widehat{\mathrm{\Lambda }}_1}\right)\widehat{\varphi }(x)`$ $`=`$ $`(i\widehat{\delta }_{\widehat{\mathrm{\Lambda }}_1}\widehat{\mathrm{\Lambda }}_2[\omega _a]i\widehat{\delta }_{\widehat{\mathrm{\Lambda }}_2}\widehat{\mathrm{\Lambda }}_1[\omega _a]+[\widehat{\mathrm{\Lambda }}_1[\omega _a]\stackrel{}{,}\widehat{\mathrm{\Lambda }}_2[\omega _a])\widehat{\varphi }(x)`$
$`=`$ $`\widehat{\mathrm{\Lambda }}_{\widehat{\mathrm{\Lambda }_1\times \mathrm{\Lambda }_2}}\widehat{\varphi }(x).`$
One finds as usual
$`[\mathrm{\Lambda }_1,\mathrm{\Lambda }_2]=i\mathrm{\Lambda }_{\mathrm{\Lambda }_1\times \mathrm{\Lambda }_2}`$ (11)
in the zeroth order in $`\theta `$ and
$`i\delta _{\mathrm{\Lambda }_1}\mathrm{\Lambda }_2^{(1)}i\delta _{\mathrm{\Lambda }_2}\mathrm{\Lambda }_1^{(1)}+i\theta ^{ab}\{_a\mathrm{\Lambda }_1,_b\mathrm{\Lambda }_2\}+[\mathrm{\Lambda }_1,\mathrm{\Lambda }_2^{(1)}][\mathrm{\Lambda }_2,\mathrm{\Lambda }_1^{(1)}]=\mathrm{\Lambda }_{\mathrm{\Lambda }_1\times \mathrm{\Lambda }_2}^{(1)}`$ (12)
in the leading order in $`\theta `$. A solution to this consistency equation is
$`\mathrm{\Lambda }_1^{(1)}={\displaystyle \frac{1}{4}}\theta ^{ab}\{_a\mathrm{\Lambda }_1,\omega _b\}`$ (13)
and analogously for $`\mathrm{\Lambda }_2^{(1)}`$ and $`\mathrm{\Lambda }_{\mathrm{\Lambda }_1\times \mathrm{\Lambda }_2}^{(1)}`$ and where we have used: $`\theta ^{ab}=\theta ^{\mu \nu }e_\mu ^ae_\nu ^b`$ and $`_a=e_a^\mu _\mu `$. Note that the consistency condition is derived in the leading order in $`\theta `$ and $`\xi (x)`$. However since $`\xi (x)`$ is itself proportional to $`\theta `$, the relevant part of the volume-preserving diffeomorphism transformation is trivial. In other words, the terms proportional to $`\xi (x)`$ can be dropped in equation (13) and it actually determines $`\lambda ^{(1)}`$ which is the first non-trivial term in the Seiberg-Witten map for the so(3,1) gauge transformation.
The consistency condition for the spin connection is given by
$`\delta _\mathrm{\Lambda }\omega _a^{(1)}`$ $`=`$ $`_a\mathrm{\Lambda }^{(1)}{\displaystyle \frac{1}{2}}\theta ^{bc}\left(_b\lambda _c\omega _a_b\omega _a_c\mathrm{\Lambda }\right)+[\mathrm{\Lambda },\omega _a^{(1)}]+i[\mathrm{\Lambda }^{(1)},\omega _a].`$ (14)
A solution is
$`\omega _a^{(1)}={\displaystyle \frac{1}{4}}\theta ^{bc}\{\omega _b,_c\omega _a+F_{ca}\}`$ (15)
One thus has $`\widehat{E}_a^\mu =e_a^\mu `$ to all orders in $`\theta `$ and $`\widehat{\omega }_a=\omega _a+\omega _a^{(1)}+๐ช(\theta ^2)`$.
The Seiberg-Witten map for the field strength is given by $`\widehat{F}_{ab}=F_{ab}+F_{ab}^{(1)}+๐ช(\theta ^2)`$, where
$`F_{ab}^{(1)}={\displaystyle \frac{1}{2}}\theta ^{cd}\{F_{ac},F_{bd}\}{\displaystyle \frac{1}{4}}\theta ^{cd}\{\omega _c,(_d+D_d)F_{ab}\},`$ (16)
where $`D_a=A_a=ie_a^\mu p_\mu +\frac{i}{2}\omega _a^{bc}\mathrm{\Sigma }_{bc}`$. is the commutative covariant derivative.
The commutative field strength $`F_{ab}`$ contains the Riemann tensor $`R_{ab}^{cd}`$ as well as a torsion $`T_{ab}^c`$:
$`F_{ab}=i[D_a,D_b]={\displaystyle \frac{1}{2}}R_{ab}^{cd}\mathrm{\Sigma }_{cd}+T_{ab}^cD_c`$ (17)
with $`R_{ab}=\frac{1}{2}R_{ab}^{cd}\mathrm{\Sigma }_{cd}`$ and $`T_{ab}^c=(D_ae_b^\nu D_be_a^\nu )e_\nu ^c`$. The commutative covariant derivative $`D_a`$ is torsion free ($`T_{ab}^c=0`$) and compatible with a metric: $`e_\mu ^aD_ae_\nu ^b=0`$. We now have all the required tools to consider actions that are invariant under general coordinate transformations.
## 3 Action for noncommutative General Relativity
The Seiberg-Witten map for the Riemann tensor $`R_{ab}`$ which is the field strength tensor corresponding to a local noncommutative Lorentz transformation can be read from equations (16) and (17) where the classical torsion is being set to zero:
$`\widehat{R}_{ab}=R_{ab}+R_{ab}^{(1)}+๐ช(\theta ^2),`$ (18)
with
$`R_{ab}^{(1)}={\displaystyle \frac{1}{2}}\theta ^{cd}\{R_{ac},R_{bd}\}{\displaystyle \frac{1}{4}}\theta ^{cd}\{\omega _c,(_d+D_d)R_{ab}\}.`$ (19)
The noncommutative Riemann tensor is then given by
$`\widehat{R}_{ab}(\widehat{x})={\displaystyle \frac{1}{2}}\widehat{R}_{ab}^{cd}(\widehat{x})\mathrm{\Sigma }_{cd},`$ (20)
from which we can determine the corresponding noncommutative Ricci tensor, $`\widehat{R}_{ab}^{bd}`$, and a Ricci scalar $`\widehat{R}=\widehat{R}_{ab}^{ab}`$ in terms of the classical fields using the above Seiberg-Witten map.
The noncommutative action is then given by
$`S`$ $`=`$ $`{\displaystyle d^4x\frac{1}{2\kappa ^2}\widehat{R}(\widehat{x})}={\displaystyle d^4x\frac{1}{2\kappa ^2}\left(R(x)+R^{(1)}(x)\right)}+๐ช(\theta ^2).`$ (21)
In the second line we have made use of the Weyl quantization procedure which allows to replace the noncommutative variables by commuting ones by expanding the noncommutative fields using the Seiberg-Witten maps. The only correction to leading order in $`\theta `$ comes from the Seiberg-Witten map for the so(3,1) gauge field. It is easy to verify that this action is Hermitian and invariant under unimodular coordinate transformations and local Lorentz transformations. This noncommutative general relativity theory is, by construction, torsion free. In the leading order in the expansion in $`\theta `$ we can use the classical relations:
$`\omega _\mu ^{ab}(x)={\displaystyle \frac{1}{2}}e_\mu ^c(x)\left(\mathrm{\Omega }_c^{ab}(x)\mathrm{\Omega }_c^{ba}(x)\mathrm{\Omega }_c^{ab}(x)\right)`$ (22)
with
$`\mathrm{\Omega }_c^{ab}(x)=e_\mu ^a(x)e_\nu ^b(x)\left(^\mu e_c^\nu (x)^\nu e_c^\mu (x)\right).`$ (23)
The equation (21) represents an action for the noncommutative version of the unimodular theory of gravitation. The unimodular theory is known to be classically equivalent to Einsteinโs General Relativity with a cosmological constant. Indeed, we can rewrite the action (21) in the form of an Einstein-Hilbert action by introducing a Lagrange multiplier $`\mathrm{\Lambda }`$ which appears to be an arbitrary integration constant:
$$S=d^4x\left(\frac{1}{2\kappa ^2}\text{det}(e_\mu ^a(x))\left(R(x)+R^{(1)}(x)\right)+\mathrm{\Lambda }\left(\text{det}(e_\mu ^a(x))1\right)+๐ช(\theta ^2)\right).$$
(24)
In deriving (24) we have used:
$`\text{det}_{}(e_\mu ^a(x))\stackrel{\mathrm{def}}{=}{\displaystyle \frac{1}{4!}}ฯต^{\mu \nu \rho \sigma }ฯต_{abcd}e_\mu ^a(x)e_\nu ^b(x)e_\rho ^c(x)e_\sigma ^d(x)=\text{det}(e_\mu ^a(x))+๐ช(\theta ^2),`$ (25)
where $``$ is the star product.
Let us now consider the weak field approximation of the noncommutative action (24). Although the theory defined by the action (24) does not admit flat spacetime as a background solution, we can still locally (for regions with volume $`V<<1/\mathrm{\Lambda }`$) expand the tetrad around flat spacetime:
$`e_a^\mu (x)=\eta _a^\mu {\displaystyle \frac{1}{2}}h_a^\mu (x).`$ (26)
Here $`h_a^\mu (x)`$ is a weak gravitational field subject to the traceless condition ($`\text{det}(e_\mu ^a(x))=1`$ follows from (24)). The noncommutative correction to Einsteinโs action in the weak field limit reads
$`{\displaystyle \frac{1}{8}}\theta ^{ln}\stackrel{~}{R}_{kl}^{ab}\stackrel{~}{R}_{mn}^{cd}d_{abcd}^{km}{\displaystyle \frac{1}{16}}\theta ^{ln}\stackrel{~}{R}_{ln}^{ab}\stackrel{~}{R}_{km}^{cd}d_{abcd}^{km}`$ (27)
where $`d_{abcd}^{km}`$ are the structure constants defined by $`d_{abcd}^{km}\mathrm{\Sigma }_{km}=2\{\mathrm{\Sigma }_{ab},\mathrm{\Sigma }_{cd}\}`$. Notice that
$`\stackrel{~}{R}_{km}^{ab}={\displaystyle \frac{1}{2}}_k(^ah_m^b+{\displaystyle \frac{1}{2}}^bh_m^a)_m(^ah_k^b^bh_k^a)`$ (28)
is the leading order of the weak field approximation of $`R_{ab}^{cd}`$. This modification of the linearized noncommutative action implies a noncommutative correction of the equations of motion for the weak gravitational field:
$`R_b^a{\displaystyle \frac{1}{2}}\eta _b^aR={\displaystyle \frac{1}{8}}\theta ^{an}\left(^r_k\stackrel{~}{R}_{mn}^{cd}+^r_n\stackrel{~}{R}_{km}^{cd}\right)d_{rbcd}^{km}{\displaystyle \frac{1}{8}}\theta ^{ln}^r_l\stackrel{~}{R}_{mn}^{cd}d_{rbcd}^{am}`$ (29)
where $`R_b^a`$ is the usual Ricci tensor and $`R`$ is the corresponding Ricci scalar in the linearized approximation (we have omitted the contribution coming from the cosmological constant). This modification might have some interesting physical implications that will be studied elsewhere.
Finally we briefly discuss the relation between the tetrad formalism considered here and a second-order formalism which involves the metric tensor. The noncommutative metric tensor defined naively as $`\widehat{g}_{\mu \nu }(\widehat{x})=E_\mu ^a(\widehat{x})E_\nu ^b(\widehat{x})\eta _{ab}`$ is neither real nor symmetric. This raises a more generic question of the geometrical interpretation of the noncommutative deformation of Einsteinโs General Relativity considered above. The simple prescription to define the metric in our case is to solve the deformed equations of motion for the classical tetrads at each given order of $`\theta `$ expansion and then to determine the metric in the standard way.
## 4 Conclusions
We have constructed a theory of noncommutative General Relativity on a canonical noncommutative spacetime. The general coordinate transformations is shown to be restricted to the volume-preserving transformations. Thus the General Relativity on canonical noncommutative spacetimes is the noncommutative version of the unimodular theory of gravitation. The local Lorentz invariance is described as a noncommutative gauge theory by taking the spin-connection field in the enveloping algebra.
The action for noncommutative General Relativity was constructed and the expansion in first order of the noncommutative parameter $`\theta `$ has been calculated. We derived the noncommutative correction to the equations of motion of the weak gravitational field. An interesting question is whether the effects of spacetime noncommutativity coming from the noncommutative modifications of gravity are stronger than the ones coming from the modifications of the interactions of the standard model .
It will also be interesting to consider classical solutions of the noncommutative action defined in this work. The minimal length introduced in our version of General Relativity could have interesting consequences for the horizon and the singularity of black holes. It will also be worth studying cosmological implications and the quantization of the noncommutative action.
### Acknowledgments
The work of X. C. was supported in part by a scholarship of the Universitรฉ libre de Bruxelles. The work of A. K. was partially supported by the GRDF grant 3305.
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# Correlated Band Structure and the Ground-State Phase Diagram in High-TC Cuprates
## Abstract
We review results obtained with a recently proposed variational cluster approach (VCA) for the competition between d-wave superconductivity (dSC) and antiferromagnetism (AF) in the high-T<sub>C</sub> cuprates. Comparing the single-particle spectra of a two-dimensional Hubbard model with quantum Monte-Carlo (QMC) and experimental data, we verify that the VCA correctly treats the low-energy excitations. The cluster calculations reproduce the overall ground-state phase diagram of the high-temperature superconductors both for electron- and hole-doping. In particular, they include salient features such as the enhanced robustness of the AF state in case of electron doping. For electron- but also for hole-doping, we clearly identify a tendency to phase separation into a mixed AF-dSC phase at low and a pure dSC-phase at high doping.
, , ,
The central issue in the field of high-temperature superconductivity (HTSC) is the connection of the microscopic interactions at the level of electrons and ions, which are at high energy and temperature $`T`$, with the โemerging phenomenaโ at $`T=0`$, i.e. competing and nearly degenerate orders โ antiferromagnetism (AF), d-wave superconductivity (dSC), heterogeneous phases, etc. We will not go into a lengthy discussion of what interactions should be retained at the electron-ion level. But, when choosing the two-dimensional $`(2D)`$ one-band Hubbard model, i.e.
$$H=\underset{i,j}{}t_{i,j}c_i^{}c_j+U\underset{i}{}n_{i,}n_{i,},$$
(1)
where $`t_{i,j}`$ denote hopping matrix elements, $`n_{i,}`$ the density at site $`i`$ with spin โ$``$โ and $`U`$ the local Coulomb repulsion, one has introduced gross simplifications, leaving out other orbital (e.g. $`p`$) degrees of freedom, long-range Coulomb interaction, electron-phonon coupling, etc . Nevertheless, this model choice appears to be legitimate, last not least in view of the amazing agreement achieved between numerical simulations and experimental results for the normal-state properties of the cuprates .
At low temperatures different orders appear, which are not separated by distinct energy scales but compete with each other. What is required is a kind of โmagnifying lensโ which allows to resolve these competing orders. Ideally, one should employ a systematic renormalization-group approach to integrate out the irrelevant degrees of freedom and, thereby, to correctly bridge high to low energies and eventually to go to $`T=0`$. For the strong-correlation case, realized in the HTSC, how to do this is, however, by no means obvious. In this context, cluster techniques provide an alternative way to systematically approach the infinite-size (and, thereby, low-energy) limit.
Here, we review progress obtained with the variational cluster approach (VCA), which was proposed and used by Potthoff et. al. . This approach provides a rather general and controlled way to go to the infinite-sized lattice fermion system at low temperatures and at $`T=0`$, in particular. The ground-state phase diagram of the $`2D`$ one-band Hubbard model was calculated within VCA by Sรฉnรฉchal et. al. and, independently, by two of us . There are certain technical differences, which we discuss below, but the โupshotโ of the two works is as follows: For the cluster sizes used in the VCA, the $`T=0`$ phase diagram of the Hubbard model (1), with hopping terms up to third-nearest neighbors, correctly reproduces salient features of the HTSC, such as the AF and dSC ground states in doping ranges, which are qualitatively in agreement with electron- and hole-doped cuprates.
The VCA is based on the self-energy-functional approach (SFA) . The SFA provides a variational scheme to use dynamical information from an exactly solvable โreference systemโ (for example an isolated cluster) to approximate the physics of a system in the thermodynamic limit. For a system with Hamiltonian $`H=H_0(๐)+H_1(๐ผ)`$ and one-particle and interaction parameters $`๐`$ and $`๐ผ`$, the grand potential is written as a functional of the self-energy $`๐บ`$ as
$$\mathrm{\Omega }_{๐,๐ผ}[๐บ]=F_๐ผ[๐บ]+\text{Tr}\mathrm{ln}(๐ฎ_{0,๐}^1๐บ)^1,$$
(2)
with the stationary property $`\delta \mathrm{\Omega }_{๐,๐ผ}[๐บ_{\mathrm{phys}}]=0`$ for the physical self-energy. Here, $`๐ฎ_{0,๐}=(\omega +\mu ๐)^1`$ is the free Greenโs function at frequency $`\omega `$, and $`\mu `$ is the chemical potential. $`F_๐ผ[๐บ]`$ is the Legendre transform of the Luttinger-Ward functional and determines the fully interacting Greenโs function via $`๐ฎ=\delta F_๐ผ[๐บ]/\delta ๐บ`$. It is important to note that $`F_๐ผ[๐บ]`$ is a universal functional: The functional dependence is only determined by the interaction parameters $`๐ผ`$ (for example, the Hubbard interaction in Eq. (1)). Therefore, the functional $`F_๐ผ[๐บ]`$ is the same as the functional for a problem which is โsimplerโ and solvable, i.e. for a Hamiltonian $`H^{}=H_0(๐^{})+H_1(๐ผ)`$ with the same interaction part but a one-particle part that makes it exactly solvable. The stationary solutions are obtained (and this is the approximation) within the subspace of self-energies $`๐บ=๐บ(๐^{})`$ of that simpler solvable problem that is spanned by varying $`๐^{}`$. If one takes a single site and connects it to continuous (non-interacting) bath degrees of freedom (another $`H^{}`$ choice), one recovers the dynamical mean-field theory (DMFT) or, for a cluster of sites connected to a bath, a cluster variant of DMFT .
In the VCA, considered in the following discourse, $`H^{}`$ is build up of disconnected clusters, which have their inter-cluster hopping terms removed. In our $`T=0`$ approach, the isolated cluster is solved by exact diagonalization. Its Hamiltonian $`H^{}`$ includes additional symmetry-breaking โWeissโ fields to allow for long-range order. The VCA solution is finally obtained as a stationary point determined by $`\mathrm{\Omega }_{๐,๐ผ}[\mathrm{\Sigma }(๐^{})]/๐^{}=0`$. Is this a controlled route to a $`(T=0)`$ infinite-size approach? To answer this, consider a few โtestsโ:
(i) The VCA correctly reproduces long-range AF order in $`2D`$ and the absence of this order in $`1D`$ . This non-trivial test implies that the VCA goes well beyond ordinary mean-field theory.
(ii) An advantage, compared to variational schemes based on wave functions , is that the VCA quite naturally gives the one-particle Greenโs function $`๐ฎ`$. Fig. 1 compares the spectral function $`A(๐,\omega )\text{Im}G(๐,\omega )`$ of the VCA for the $`2D`$ Hubbard model at $`U=8t`$, half-filling and $`T=0`$ with corresponding low-temperature QMC data for an isolated $`8\times 8`$ cluster. One can clearly see that the VCA, with the lattice covered by $`\sqrt{10}\times \sqrt{10}`$ clusters, correctly reproduces coherent and incoherent โbandsโ (known from ARPES data ). In particular, the non-trivial proliferation of AF spin correlations from cluster to cluster, which builds up the coherent quasi-particle โbandโ is obviously correctly embedded in the VCA . Similar calculations have been performed for the spectral function $`A(๐,\omega )`$ of the hole- and electron-doped Hubbard model . The characteristically different doping dependencies give rise to different Fermi-surface evolutions upon doping. Furthermore, the single-particle excitations provide insight into the characteristic differences in the ground-state phase diagram for hole- and electron-doping .
(iii) To test the stability of the homogeneous phases with respect to phase separation (PS), we consider a reference system $`H^{}`$ of isolated $`2\times 2`$ clusters where, in addition to the two symmetry-breaking terms (Weiss fields) $`H_{\mathrm{AF}}^{}`$ and $`H_{\mathrm{SC}}^{}`$, a term $`H_{local}^{}`$ is optimized within the variational procedure . $`H_{local}^{}=\epsilon _{i\sigma }n_{i\sigma }`$ describes a shift $`\epsilon `$ of the chemical potential in the cluster with respect to the physical chemical potential $`\mu `$. The use of the additional variational parameter $`\epsilon `$ is required in order to have a consistent treatment of the particle density. The optimization of $`\epsilon `$ has to be done simultaneously with the optimization of the parameters $`h_{\mathrm{AF}}`$ and $`h_{\mathrm{SC}}`$, namely the staggered magnetic field in the term $`H_{\mathrm{AF}}^{}`$ and the nearest-neighbor d-wave pairing field $`h_{\mathrm{SC}}`$ in the term $`H_{\mathrm{SC}}^{}`$. Notice that the Weiss fields $`h_{\mathrm{AF}}`$ and $`h_{\mathrm{SC}}`$ are different from the corresponding order parameters $`m`$ and $`\mathrm{\Delta }`$ plotted in Figs. 2 and 3. Quite generally, however, a nonvanishing stationary value for the Weiss fields produces a nonvanishing order parameter, although the latter can be much smaller.
The phase diagram for the Hubbard model with $`U=8t`$ and next-nearest-neighbor hopping $`t_{n.n.n}=0.3t`$, obtained with our calculation, is plotted in Fig. 2 for the hole-doped and in Fig. 3 for the electron-doped case. In the upper part of each figure, we display the AF $`(m)`$ and dSC $`(\mathrm{\Delta })`$ order parameters as a function of doping $`x`$. In the lower part of the figures, the chemical potential $`\mu `$ is plotted as a function of $`x`$.
Let us discuss hole doping first (see Fig. 2). For dopings $`x`$ below a critical value $`x_1`$ we find a homogeneous symmetry-broken state in which both, the AF as well as the dSC order parameter $`m`$ and $`\mathrm{\Delta }`$ are non-zero. This corresponds to a phase AF+SC where AF and dSC order microscopically and coherently coexist. A homogeneous state with pure dSC ($`m=0`$ and $`\mathrm{\Delta }>0`$) is obtained for dopings $`x>x_2`$.
Fig. 2 also shows $`\mathrm{\Delta }`$ and $`\mu `$ for the homogeneous AF+SC and SC phases in the range $`x_1<x<x_2`$. Here, however, these phases are thermodynamically unstable. For dopings $`x`$ with $`x_1<x<x_2`$, macroscopic phase separation between the two phases occurs. In practice, doping-dependencies are calculated by varying the chemical potential $`\mu `$. Following up the grand potentials for the two homogeneous phases as functions of $`\mu `$, i.e. $`\mathrm{\Omega }_{\mathrm{AF}+\mathrm{SC}}(\mu )`$ and $`\mathrm{\Omega }_{\mathrm{SC}}(\mu )`$, it is found (see Ref. for details) that there is a crossing at a critical chemical potential $`\mu =\mu _c`$ (at this point the AF order parameter $`m`$ is still nonzero). Thus, the transition is first order as a function of $`\mu `$. At the transition point $`\mu _c`$, the dopings corresponding to the AF+SC and to the SC phase, $`x_{\mathrm{AF}+\mathrm{SC}}`$ and $`x_{\mathrm{SC}}`$, are different. Consequently, there is a jump $`\mathrm{\Delta }xx_{\mathrm{AF}+\mathrm{SC}}x_{\mathrm{SC}}`$ at $`\mu _c`$, indicating phase separation between a weakly doped AF+SC and a higher doped SC phase.
Due to the inclusion of $`H_{\mathrm{local}}^{}`$, $`\mu _c`$ can equivalently be obtained by a Maxwell construction. This is shown in Fig. 2 where, in the lower panel, $`\mu `$ is plotted as a function of $`x`$. Here, phase separation is signaled by the fact that the $`\mu (x)`$ is not a monotonous function. The Maxwell construction shown in the figure then identifies the two dopings $`x_1`$ and $`x_2`$ into which the system tends to phase separate, as well as the chemical potential $`\mu _c`$ in the phase-separated state. In Fig. 2, $`\mu ^{}`$ is the point where the slope of $`\mu (x)`$ changes sign. For $`\mu <\mu ^{}`$ the AF+SC solution ceases to exist
Let us now discuss the electron-doped case. While the phase diagrams in Figs. 2 and 3 are qualitatively similar, the phase in which long-range AF order is realized is spreading to significantly larger doping values in the electron-doped case, in overall agreement with the experimental situation. Another important difference concerns the energy scale for phase separation, i.e. $`\mathrm{\Delta }\mu (\mu ^{}\mu _c)`$. As one can see from the comparison between Figs. 2 and 3, $`\mathrm{\Delta }\mu `$ is an order of magnitude larger in the hole-doped case. In Ref. it is argued that this can explain the different pseudogap and SC transition scales in hole- and electron-doped materials. This may give support to theories which are based on the idea that fluctuations of the competing phases, or of the related order parameters, are responsible for the pseudogap phenomenon.
In order to resolve the relevant small energy scale, it is necessary to evaluate $`\mathrm{\Omega }`$ as well as its stationary points with high accuracy. Furthermore, the inclusion of the chemical potential shift term $`H_{local}^{}`$ considerably complicates the variational optimization. For the rather small clusters of size $`2\times 2`$ used here, the reference system can be treated by full diagonalization and the frequency integrals, which are implicit in Eq. (2) , can be carried out by means of a sum over the negative poles of the Greenโs functions. Sรฉnรฉchal et. al. have considered clusters up to 10 sites and report similar results for $`x0`$ but, without the inclusion of $`H_{local}^{}`$, one cannot reliably test the stability of the homogeneous phases against phase separation. Note that the transition from the AF+SC to the SC phase may appear as continuous as a function of $`x`$ if phase separation is not taken into account.
In conclusion, there has been substantial recent progress in relating the โhigh-energyโ physics of the Hubbard model and its variants to the low-energy physics of the competing phases AF, dSC, charge inhomogenities, etc. This progress is due to the development of quantum-cluster theories, such as the VCA discussed here but also due to cluster extensions of the DMFT, such as the dynamical cluster approximation (DCA) . Within these cluster approaches, characteristic difficulties have been encountered: The latest impressive work using the DCA by Maier et. al. performed a systematic cluster-size study of dSC in the $`2D`$ Hubbard model . In clusters large enough (up to 26 sites) converged results point to a finite-$`T`$ instability to dSC. Because of the QMC minus-sign problem, however, results were limited to $`U=4t`$, where the typical energy separation in $`U`$ and the magnetic energy scale of the HTSC is not yet achieved. On the other hand, the VCA studies reviewed here, are clearly not yet converged with respect to the cluster size, as one can read off from Fig. 1 in Ref . Cluster convergence is, at least in principle, also possible within the $`VCA`$. With increasing cluster size longer-ranged correlations are included exactly. This, however, necessarily implies to use stochastic (QMC) methods as solvers for the cluster reference system.
One of us (WH) would like to acknowledge the hospitality of the Kavli Institute for Theoretical Physics in Santa Barbara, where part of this work was finished (supported by Nat. Sc. Found. Grant No. PHY99-0794). We would like to thank D. J. Scalapino for many useful discussions. This work was also supported by the DFG Forschergruppe 538, by the KONWIHR supercomputing network in Bavaria, and by the Doctoral Scholarship Program of the Austrian Academy of Sciences.
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# Ramanujanโs Perimeter of an Ellipse
## 1 Introduction
Let $`a`$ and $`b`$ be the semi-major and semi-minor axes of an ellipse with perimeter $`p`$ and whose eccentricity is $`k`$. The final sentence of Ramanujanโs famous paper *Modular Equations and Approximations to $`\pi `$,* , says:
> *โ The following approximation for* $`p`$ \[was\] obtained empirically:
>
> $$\overline{)p=\pi \left\{(a+b)+\frac{3(ab)^2}{10(a+b)+\sqrt{a^2+14ab+b^2}}+\epsilon \right\}}$$
> (1.1)
> *where* $`\epsilon `$ is about $`{\displaystyle \frac{3ak^{20}}{68719476736}}.`$
Ramanujan never explained his โempiricalโ method of obtaining this approximation, nor ever subsequently returned to this approximation, neither in his published work, nor in his Notebooks . Indeed, although the Notebooks does contain the above approximation (see Entry 3 of Chapter XVIII) the statement there does not even mention his asymptotic error estimate stated above.
Twenty years later Watson claimed to have proven that Ramanujanโs approximation is *in defect*, but he never published his proof.
In 1978, we established the following optimal version of Ramanujanโs approximation:
###### THEOREM 1.
(Ramanujanโs Approximation Theorem) Ramanujanโs approximative perimeter
$$\overline{)p_R:=\pi \left\{(a+b)+\frac{3(ab)^2}{10(a+b)+\sqrt{a^2+14ab+b^2}}\right\}}$$
(1.2)
underestimates the true perimeter, $`\mathrm{p}`$, by
$$\overline{)ฯต:=\pi (a+b)\theta (\lambda )\lambda ^{10},}$$
(1.3)
where
$$\lambda :=\frac{ab}{a+b},$$
(1.4)
and where the function $`\theta (\lambda )`$ grows monotonically in $`0\lambda 1`$ while at the same time it satisfies the optimal inequalities
$$\overline{)\frac{3}{2^{17}}<\theta (\lambda )\frac{14}{11}\left(\frac{22}{7}\pi \right)}$$
(1.5)
$`\mathrm{}`$
Please take note of the striking form of the sharp upper bound since it involves the number $`\left({\displaystyle \frac{22}{7}}\pi \right)`$ which measures *the accuracy of Archimedesโ famous approximation, $`{\displaystyle \frac{22}{7}},`$ to the transcendental number $`\pi `$!*
###### COROLLARY 1.
The error in defect, $`ฯต`$, as a function of $`\lambda `$, grows monotonically for $`0\lambda 1.`$
$`\mathrm{}`$
###### COROLLARY 2.
The error in defect, $`ฯต`$, as a function of the eccentricity, $`e`$, is given by
$$\overline{)ฯต(e):=a\left\{\delta (e)\left(\frac{2}{1+\sqrt{1e^2}}\right)^{19}\right\}e^{20}}.$$
(1.6)
Moreover, $`ฯต(e)`$ grows monotonically with $`e`$, $`0e1`$, while $`\delta (e)`$ satisfies the optimal inequalities
$$\overline{)\frac{3\pi }{68719476736}<\delta (e)\frac{{\displaystyle \frac{7}{11}}\left({\displaystyle \frac{22}{7}}\pi \right)}{2^{18}}}$$
(1.7)
$`\mathrm{}`$
This Corollary 2 explains the significance of Ramanujanโs own error estimate in (1.1). The latter is an asymptotic *lower bound* for $`ฯต(e)`$ but it is not the optimal one. That is given in (1.7).
## 2 Later History
We sent an (updated) copy of our 1978 preprint to Professor Bruce Berndt in 1988 and he subsequently quoted its conclusions in his edition of Volume 3 of the Notebooks (see p. 150 ). However the details of our proofs have never been published and so we have decided to present them in this paper.
Berndtโs discussion of Ramanujanโs approximation includes Almkvistโs very plausible suggestion that Ramanujanโs โempirical processโ was to develop a *continued fraction expansion* of Ivoryโs infinite series for the perimeter () as well as a proof, due independently to Almkvist and Askey, of our fundamental lemma (see ยง3). However, their proof is different from ours.
The most recent work on the subject has been carried out by R. Barnard, K. Pearce, and K. Richards in and was published in the year $`2000`$. They also prove the major conclusion in our fundamental lemma, but their methods too are quite different from ours.
## 3 Fundamental Lemma
###### THEOREM 2.
(Fundamental Lemma) Define the functions $`๐(x)`$ and $`๐(x)`$ and the coefficients $`A_n`$ and $`B_n`$ by:
$`๐(x)`$ $`:=1+{\displaystyle \frac{3x}{10+\sqrt{43x}}}:=1+A_1x+A_2x^2+\mathrm{}`$ (3.1)
$`๐(x)`$ $`:={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left\{{\displaystyle \frac{1}{2n1}}{\displaystyle \frac{1}{4^n}}\left({\displaystyle \genfrac{}{}{0pt}{}{2n}{n}}\right)\right\}^2x^n:=1+B_1x+B_2x^2+\mathrm{}.`$ (3.2)
Then:
$`A_1=B_1,A_2=B_2,A_3=B_3,A_4=B_4`$ (3.3)
$`A_5<B_5,A_6<B_6,\mathrm{},A_n<B_n,\mathrm{}`$ (3.4)
where the strict inequalities in (3.4) are valid for all $`n5.`$
###### Proof.
First we prove $`(\mathbf{3.3})`$. We read this off directly from the numerical values of the expansion:
$`A_1`$ $`=B_1={\displaystyle \frac{1}{4}}`$
$`A_2`$ $`=B_2={\displaystyle \frac{1}{16}}`$
$`A_3`$ $`=B_3={\displaystyle \frac{1}{64}}`$
$`A_4`$ $`=B_4={\displaystyle \frac{25}{4096}}.`$
Now we prove $`(\mathbf{3.4})`$ For $`A_5`$, $`B_5`$, $`A_6`$, and $`B_6`$ we verify $`(3.4)`$ directly from their explicit numerical values. Namely,
$`A_5`$ $`={\displaystyle \frac{47\frac{1}{2}}{2^{14}}},B_5={\displaystyle \frac{49}{2^{14}}},A_5B_5={\displaystyle \frac{\frac{3}{2}}{2^{14}}}<0`$
$`A_6`$ $`={\displaystyle \frac{803}{2^{21}}},B_6={\displaystyle \frac{882}{2^{21}}},A_6B_6={\displaystyle \frac{79}{2^{21}}}<0.`$
*Therefore it is sufficient to prove*
$$A_n<B_n$$
(3.5)
*for all*
$$n7.$$
(3.6)
Now the *explicit* formula for $`A_n`$ is
$$A_n=a_{n1}+a_{n2}+a_{n3}+\mathrm{}+a_1+a_0$$
(3.7)
where
$$\overline{)\begin{array}{ccc}\hfill a_{n1}& :=\frac{1}{2n3}\frac{1}{16^n}\left(\genfrac{}{}{0pt}{}{2n2}{n1}\right)3^{n1}\hfill & \\ \hfill a_{n2}& :=\frac{1}{2n5}\frac{1}{16^{n1}}\left(\genfrac{}{}{0pt}{}{2n2}{n1}\right)3^{n2}\left(\frac{1}{2^5}\right)\hfill & \\ \hfill .& .\hfill & \\ \hfill .& .\hfill & \\ \hfill .& .\hfill & \\ \hfill a_1& :=\frac{1}{211}\frac{1}{16^2}\left(\genfrac{}{}{0pt}{}{2}{1}\right)3^{n2}\left(\frac{1}{2^5}\right)^{n2}\hfill & \\ \hfill a_0& :=\frac{4}{16}\left(\frac{1}{2^5}\right)^{n1}.\hfill & \end{array}}$$
(3.8)
Next we write
$$A_n=a_{n1}\left(1+\frac{a_{n2}}{a_{n1}}+\frac{a_{n3}}{a_{n1}}+\frac{a_{n4}}{a_{n1}}+\mathrm{}+\frac{a_1}{a_{n1}}+\frac{a_0}{a_{n1}}\right)$$
(3.9)
and assert:
###### CLAIM 1.
*The ratios $`{\displaystyle \frac{a_{nk1}}{a_{nk}}}`$ decrease monotonically in absolute value as $`k`$ increases from $`k=1`$ to $`k=n1.`$*
###### Proof.
For $`k=1,\mathrm{},n2,`$
$`\left|{\displaystyle \frac{a_{nk1}}{a_{nk}}}\right|`$ $`=\left(1+{\displaystyle \frac{2}{2n2k3}}\right)\left(1+{\displaystyle \frac{1}{4n4k2}}\right){\displaystyle \frac{1}{12}}`$
$`{\displaystyle \frac{1}{6}}(\text{which is the worst case and occurs when }k=n2\text{)}`$
$`<1`$
For $`k=n1`$,
$$\left|\frac{a_0}{a_1}\right|=\frac{1}{3}<1.$$
This completes the proof. โ
###### CLAIM 2.
*The ratios $`{\displaystyle \frac{a_{nk1}}{a_{nk}}}`$alternate in sign.*
###### Proof.
This is a consequence of the definition of the $`a_k.`$
By CLAIM 1. and CLAIM 2. we can write $`(3.9)`$ in the form
$`A_n`$ $`=a_{n1}(1\text{something positive and smaller than }1)`$
$`<a_{n1}.`$
Therefore, to prove $`(3.8)`$ for $`n7`$, *it suffices to prove*
$$a_{n1}<B_n$$
(3.10)
*for all* $`n7`$.
By $`(3.8)`$ and the definition of $`B_n`$, this last afirmation is equivalent to proving
$$\frac{1}{2n3}\frac{1}{16^n}\left(\genfrac{}{}{0pt}{}{2n2}{n1}\right)3^{n1}<\left\{\frac{1}{2n1}\frac{1}{4^n}\left(\genfrac{}{}{0pt}{}{2n}{n}\right)\right\}^2,$$
which, after some algebra, *reduces to proving the implication*
$$n7\frac{{\displaystyle \frac{n}{2}}{\displaystyle \frac{2n1}{2n3}}}{\left({\displaystyle \genfrac{}{}{0pt}{}{2n}{n}}\right)}3^{n1}<1.$$
If we define for all integers $`n7`$
$$f(n):=\frac{{\displaystyle \frac{n}{2}}{\displaystyle \frac{2n1}{2n3}}}{\left({\displaystyle \genfrac{}{}{0pt}{}{2n}{n}}\right)}3^{n1}$$
(3.11)
then the affirmation $`(3.10)`$ turns out to be *equivalent to*
$$n7f(n)<1$$
(3.12)
This latter affirmation is a consequence of the following two conditions:
1. $`f(7)<1.`$
2. $`f(7)>f(8)>f(9)>\mathrm{}>f(k)>f(k+1)>\mathrm{}`$
Proof of 1. By direct numerical computation
$$f(7)=\frac{1701}{1936}<1$$
$`\mathrm{}`$
Proof of 2. We must show
$$k7f(k)>f(k+1).$$
If we define
$$g(k):=\frac{f(k)}{f(k+1)},$$
(3.13)
then we must show
$$k7g(k)>1.$$
(3.14)
Using the definition $`(3.11)`$ of $`f(n)`$ and the definition $`(3.14)`$ of $`g(n)`$, and reducing algebraically we find
$$g(k)=\frac{2k}{6k9}\left(\frac{2k1}{k+1}\right)^2,$$
and we must show that
$$k7\frac{2k}{6k9}\left(\frac{2k1}{k+1}\right)^2>1.$$
(3.15)
Define the rational function of the real variable $`x`$:
$$g(x):=\frac{2x}{6x9}\left(\frac{2x1}{x+1}\right)^2.$$
(3.16)
Then the graph of $`y=g(x)`$ has a vertical asymptote at $`x=\frac{3}{2}`$ and
$$\underset{x\frac{3}{2}^+}{lim}g(x)=+\mathrm{}.$$
(3.17)
Moreover, the derivative of $`g(x)`$ is given by:
$$g^{}(x)=\frac{2(2x^27x+1)}{x(x+1)(2x1)(2x+3)}$$
which implies that
$$g^{}(x)\{\begin{array}{cc}<0\hfill & \text{if }\frac{3}{2}<x<\frac{7+\sqrt{41}}{4}\text{ },\hfill \\ =0\hfill & \text{if }x=\frac{7+\sqrt{41}}{4},\hfill \\ >0\hfill & \text{if }x>\frac{7+\sqrt{41}}{4}.\hfill \end{array}$$
Therefore, for $`x\frac{3}{2}`$ g(x) *decreases* from โ$`+\mathrm{}`$โ at $`x=\frac{3}{2}`$ (see $`(3.17)`$) to an *absolute minimum value* (in $`\frac{3}{2}x<\mathrm{}`$)
$$g\left(\frac{7+\sqrt{41}}{4}\right)=1+\frac{37\sqrt{41}}{399+69\sqrt{41}}=1.0363895208\mathrm{}$$
and then *increases monotonically* as $`x\mathrm{}`$ to its asymptotic limit $`y=\frac{4}{3}.`$Therefore
$`g(x)`$ $`>1.03638\mathrm{}\text{for all }x>\frac{3}{2},x\frac{7+\sqrt{41}}{4}`$
$`g(x)`$ $`>1.03638\mathrm{}\text{for all integers }n2.`$
$`f(n)`$ $`>(1.03638\mathrm{})f(n+1)\text{for all integers }n2.`$
$`f(n)`$ $`>f(n+1)\text{for all integers }n2.`$
$`f(7)`$ $`>f(8)>f(9)>\mathrm{}`$
which implies that the condition 2. holds. Moreover we conclude that
$`f(n)<1\text{for all integers }n7`$
$``$ $`(3.10)\text{holds for all integers }n7`$
$``$ $`(3.5)\text{holds for all integers }n7`$
$``$ $`(3.4)\text{holds for all integers }n5`$
and this completes the proof of the Fundamental Lemma. โ
## 4 Ivoryโs Identity
In $`1796`$, J. Ivory published the following identity (in somewhat different notation):
###### THEOREM 3.
(Ivoryโs Identity) If $`0x1`$ then the following formula for $`๐(x)`$ is valid:
$$\overline{)\frac{1}{\pi }_0^\pi \sqrt{1+2\sqrt{x}\mathrm{cos}(2\varphi )+x}๐\varphi =\underset{n=0}{\overset{\mathrm{}}{}}\left\{\frac{1}{2n1}\frac{1}{4^n}\left(\genfrac{}{}{0pt}{}{2n}{n}\right)\right\}^2x^n๐(x)}$$
(4.1)
###### Proof.
We sketch his elegant proof.
$`{\displaystyle \frac{1}{\pi }}{\displaystyle _0^\pi }\sqrt{1+2\sqrt{x}\mathrm{cos}(2\varphi )+x}๐\varphi ={\displaystyle \frac{1}{\pi }}{\displaystyle _0^\pi }\sqrt{1+\sqrt{x}e^{2i\varphi }}\sqrt{1+\sqrt{x}e^{2i\varphi }}๐\varphi `$
$`={\displaystyle \frac{1}{\pi }}{\displaystyle _0^\pi }{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}\left\{{\displaystyle \frac{1}{2m1}}{\displaystyle \frac{1}{4^m}}\left({\displaystyle \genfrac{}{}{0pt}{}{2m}{m}}\right)(\sqrt{x})^me^{2\pi im\varphi }\right\}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left\{{\displaystyle \frac{1}{2n1}}{\displaystyle \frac{1}{4^n}}\left({\displaystyle \genfrac{}{}{0pt}{}{2n}{n}}\right)(\sqrt{x})^ne^{2\pi in\varphi }\right\}d\varphi `$
$`={\displaystyle \frac{1}{\pi }}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}\left\{{\displaystyle \frac{1}{2m1}}{\displaystyle \frac{1}{4^m}}\left({\displaystyle \genfrac{}{}{0pt}{}{2m}{m}}\right)(\sqrt{x})^m\right\}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left\{{\displaystyle \frac{1}{2n1}}{\displaystyle \frac{1}{4^n}}\left({\displaystyle \genfrac{}{}{0pt}{}{2n}{n}}\right)(\sqrt{x})^n\right\}{\displaystyle _0^\pi }e^{2\pi i(mn)\varphi }๐\varphi `$
$`={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left\{{\displaystyle \frac{1}{2n1}}{\displaystyle \frac{1}{4^n}}\left({\displaystyle \genfrac{}{}{0pt}{}{2n}{n}}\right)\right\}^2x^n`$
We will need the following evaluation in our investigation of the accuracy of Ramanujanโs approximation.
###### COROLLARY 1.
$$\overline{)๐(1)=\frac{4}{\pi }}$$
(4.2)
###### Proof.
By Ivoryโs identity,
$`๐(1)=`$ $`{\displaystyle \frac{1}{\pi }}{\displaystyle _0^\pi }\sqrt{1+2\sqrt{1}\mathrm{cos}(2\varphi )+1}๐\varphi `$
$`={\displaystyle \frac{1}{\pi }}{\displaystyle _0^\pi }\sqrt{2+2\mathrm{cos}(2\varphi )}๐\varphi `$
$`={\displaystyle \frac{1}{\pi }}{\displaystyle _0^\pi }\sqrt{4\mathrm{cos}^2(\varphi )}๐\varphi `$
$`={\displaystyle \frac{4}{\pi }}`$
## 5 The Accuracy Lemma
###### THEOREM 4.
(Accuracy Lemma) For $`0x1`$, the function
$$๐(x):=1+\frac{3x}{10+\sqrt{43x}}$$
(5.1)
underestimates the function
$$๐(x):=\underset{n=0}{\overset{\mathrm{}}{}}\left\{\frac{1}{2n1}\frac{1}{4^n}\left(\genfrac{}{}{0pt}{}{2n}{n}\right)\right\}^2x^n$$
(5.2)
by a discrepancy, $`\mathrm{\Delta }(x)`$ which is never more than $`\left({\displaystyle \frac{4}{\pi }}{\displaystyle \frac{14}{11}}\right)x^5`$ and which is always more than $`{\displaystyle \frac{3}{2^{17}}}x^5`$:
$$\overline{)\frac{3}{2^{17}}x^5<\mathrm{\Delta }(x)\left(\frac{4}{\pi }\frac{14}{11}\right)x^5}$$
(5.3)
Moreover, the constants $`\left({\displaystyle \frac{4}{\pi }}{\displaystyle \frac{14}{11}}\right)`$ and $`{\displaystyle \frac{3}{2^{17}}}x^5`$ are the best possible.
###### Proof.
By the definition of $`๐(x)`$ and $`๐(x)`$ given in Theorem 1., the discrepancy $`\mathrm{\Delta }(x)`$ is given by the series
$`\mathrm{\Delta }(x)`$ $`:=๐(x)๐(x)`$
$`=(B_5A_5)x^5+(B_6A_6)x^6+\mathrm{}`$
$`:=\delta _5x^5+\delta _6x^6+\mathrm{},`$
where, again by Theorem 1.,
$$\delta _k>0\text{for }k=5,6,\mathrm{}.$$
On the one hand
$`\mathrm{\Delta }(x)`$ $`=x^5(\delta _5+\delta _6x+\mathrm{})`$
$`x^5(\delta _5+\delta _6+\delta _7+\mathrm{})`$
$`=x^5\mathrm{\Delta }(1)`$
$`=x^5\{๐(1)๐(1)\}`$
$`=x^5\left({\displaystyle \frac{4}{\pi }}{\displaystyle \frac{14}{11}}\right)`$
where we used Corollary 1 of Ivoryโs identity in the last equality. Therefore
$$\mathrm{\Delta }(x)\left(\frac{4}{\pi }\frac{14}{11}\right)x^5.$$
This is half of the accuracy lemma. Moreover the constant $`\left({\displaystyle \frac{4}{\pi }}{\displaystyle \frac{14}{11}}\right)`$ is *assumed* for $`x=1`$ and thus cannot be replaced by anything smaller, i.e., it is *the best possible* constant.
On the other hand, we can write
$$\mathrm{\Delta }(x)=x^5\{\delta _5+G(x)\},$$
where
$$G(x):=\delta _6x+\delta _7x^2+\mathrm{}\{\begin{array}{cc}G(x)0\hfill & \text{for all }0x1\text{ },\hfill \\ G(x)0\hfill & \text{as }x0.\hfill \end{array}$$
This shows that
$$\mathrm{\Delta }(x)>\delta _5x^5=\frac{3}{2^{17}}x^5$$
and that
$$\underset{x0}{lim}\frac{\mathrm{\Delta }(x)}{x^5}=\frac{3}{2^{17}}.$$
This proves both the other inequality in the theorem and the *optimality* of the constant $`\delta _5={\displaystyle \frac{3}{2^{17}}},`$ i.e., that it cannot be replaced by any larger constant.
This completes the proof of the Accuracy Lemma.โ
## 6 The Accuracy of Ramanujanโs Approximation
Now we can achieve the main goal of this paper, namely to prove *Ramanujanโs Approximation Theorem*.
First we express the perimeter of an ellipse and Ramanujanโs approximative perimeter in terms of the functions $`๐(x)`$ and $`๐(x)`$.
###### THEOREM 5.
If $`p`$ is the perimeter of an ellipse with semimajor axes $`a`$ and $`b`$, and if $`p_R`$ is Ramanujanโs approximative perimeter, then:
$$\overline{)\begin{array}{ccc}\hfill p& =\pi (a+b)๐\left\{\left(\frac{ab}{a+b}\right)^2\right\}\hfill & \\ & & \\ \hfill p_R& =\pi (a+b)๐\left\{\left(\frac{ab}{a+b}\right)^2\right\}.\hfill & \end{array}}$$
(6.1)
###### Proof.
We begin with *Ivoryโs Identity* (ยง4) and in it we substitute $`x:=\left({\displaystyle \frac{ab}{a+b}}\right)^2.`$ Then the integral becomes
$`{\displaystyle \frac{1}{\pi }}{\displaystyle _0^\pi }\sqrt{1+2\sqrt{\left({\displaystyle \frac{ab}{a+b}}\right)^2}\mathrm{cos}(2\varphi )+\left({\displaystyle \frac{ab}{a+b}}\right)^2}๐\varphi `$ $`={\displaystyle \frac{4}{\pi (a+b)}}{\displaystyle _0^{\frac{\pi }{2}}}(a^2\mathrm{sin}^2\varphi +b^2\mathrm{cos}^2\varphi )๐\varphi `$
and therefore
$`๐\left\{\left({\displaystyle \frac{ab}{a+b}}\right)^2\right\}`$ $`={\displaystyle \frac{4}{\pi (a+b)}}{\displaystyle _0^{\frac{\pi }{2}}}(a^2\mathrm{sin}^2\varphi +b^2\mathrm{cos}^2\varphi )๐\varphi `$
But, it is well known (Berndt ) that *the perimeter, $`\mathrm{p}`$, of an ellipse with semiaxes $`\mathrm{a}`$ and $`\mathrm{b}`$ is given by*
$$p=4_0^{\frac{\pi }{2}}(a^2\mathrm{sin}^2\varphi +b^2\mathrm{cos}^2\varphi )๐\varphi ,$$
and thus
$$p=\pi (a+b)๐\left\{\left(\frac{ab}{a+b}\right)^2\right\}.$$
(6.2)
Moreover, some algebra shows us that
$`๐\left\{\left({\displaystyle \frac{ab}{a+b}}\right)^2\right\}`$ $`=1+{\displaystyle \frac{3\left({\displaystyle \frac{ab}{a+b}}\right)^2}{10+\sqrt{43\left({\displaystyle \frac{ab}{a+b}}\right)^2}}}`$
$`={\displaystyle \frac{1}{a+b}}\left\{(a+b)+{\displaystyle \frac{3(ab)^2}{10(a+b)+\sqrt{a^2+14ab+b^2}}}\right\}`$
and we conclude that Ramanujanโs *approximative formula, $`\mathrm{p}_\mathrm{R}`$ is given by*
$$p_R=\pi (a+b)๐\left\{\left(\frac{ab}{a+b}\right)^2\right\}.$$
(6.3)
The formula for $`p`$ above was the object of Ivoryโs original paper .
Now we complete the proof of Theorem 1.
###### Proof.
Writing
$$\lambda :=\frac{ab}{a+b},$$
and using the notation of the statement of Theorem 1. we conclude that
$`ฯต`$ $`:=\pi (a+b)\theta (\lambda )\lambda ^{10}`$
$`=\pi (a+b){\displaystyle \frac{\mathrm{\Delta }(\lambda ^2)}{\lambda ^{10}}}\lambda ^{10}`$
where
$$\theta (\lambda )\frac{\mathrm{\Delta }(\lambda ^2)}{\lambda ^{10}}=\delta _5+\delta _6\lambda ^2+\mathrm{}.$$
(6.4)
Now we apply the *Accuracy Lemma* and the proof is complete.
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# A sample of radio-loud AGN in the Sloan Digital Sky Survey
## 1 Introduction
In recent years, new radio surveys such as the National Radio Astronomy Observatory (NRAO) Very Large Array (VLA) Sky Survey (NVSS; Condon et al. 1998) and the Faint Images of the Radio Sky at Twenty centimetres (FIRST) survey \[Becker et al. 1995\] have covered substantial fractions of the sky down to milli-Jansky flux densities, at vastly higher angular resolution than their predecessors. Such radio surveys are dramatically advancing our understanding of extragalactic radio sources, by permitting detailed statistical studies to be carried out. In order to reap the full benefit of these surveys, it is necessary to optically identify the radio sources, so as to obtain spectroscopic redshifts and determine the properties of their host galaxies. The availability of new large galaxy redshift surveys, especially the 2-degree Field Galaxy Redshift Survey (2dFGRS; Colless et al. 2001) and the Sloan Digital Sky Survey (SDSS; York et al. 2000; Stoughton et al. 2002), means that optical identifications and redshifts are available for large samples of nearby radio sources and allows comprehensive statistical analyses of their host galaxy properties to be carried out.
Automated crossโcorrelation of surveys across different wavelength regimes has a long history in astronomy. It is important that this process should maximize both the completeness and the reliability of the resulting sample, so considerable care needs to be taken in choosing the parameters that determine whether objects in different catalogues are indeed associated. In the case of optical identification of radio sources, the choice of radio survey is important. This is because many radio sources are extended, with sizes from a few arcsec up to tens of arcmins, and in high angular resolution surveys different components of the same source may be resolved into distinct sources. Surveys at lower angular resolution detect most sources as single components, and also have good sensitivity to extended radio structures, but the high surface density of possible optical counterparts (over 4000 per square degree at high Galactic latitudes in the Palomar Observatory Sky Survey) limits the reliability of the optical matching.
The first radio sky surveys were carried out at very low angular resolution and detected only the brightest radio sources. The resolution of these surveys was too low to allow identification of the host galaxies without detailed radio follow-up observations of the detected sources; this was time-consuming and meant that only small samples of galaxies could be studied (see discussion in McMahon et al. 2002). The NVSS was the first radio survey of sufficiently high angular resolution (45 arcsec) to permit automated crossโcorrelation with optical surveys. Machalski & Condon \[Machalski & Condon 1999\] crossโcorrelated the NVSS with the Las Campanas Redshift Survey (LCRS; Shectman et al. 1996), identifying 1157 radioโemitting galaxies. Machalski & Godlowski \[Machalski & Godlowski 2000\] used this sample to derive the local radio luminosity function. Using farโinfrared data available for the LCRS they were also able to separate the luminosity function into a radioโloud active galactic nuclei (AGN) component, which dominates at high radio luminosities, and a lower-luminosity component due to starโforming galaxies that emit in the radio predominantly due to the synchrotron emission from supernova remnants. Similarly, Sadler et al. \[Sadler et al. 2002\] cross-correlated the NVSS with galaxies from the first data release of the 2dFGRS, defining a sample of 912 radio sources which form a basis for further detailed studies (e.g. Best 2004).
The 45 arcsec resolution of the NVSS has the advantage of being sufficiently large that $`99`$% of radio sources are contained within a single NVSS component. With the exception of a few very large sources, the NVSS is also able to detect the entirety of the radio emission. However, the poor angular resolution of the NVSS leads to significant uncertainties in cross-identifying the radio sources with their optical host galaxies and there is a trade-off between the reliability of the matched sample and its completeness. Sadler et al. \[Sadler et al. 2002\] accepted radio sources within a matching radius of 10 arcseconds from an optical galaxy, leading to a catalogue that was $`90`$% complete, but in which 5โ10% of the matches are expected to be false identifications.
Samples with much higher reliability can be derived using the FIRST catalogue, due to its superior angular resolution ($`5`$ arcsec). Iveziฤ et al. \[Iveziฤ et al. 2002\] crossโcorrelated the FIRST survey with the SDSS imaging sample. Under the assumption that all true identifications of point radio sources would have radioโoptical positional offsets of less than 3 arcsec, they concluded that the optimal matching radius for crossโcorrelation was 1.5 arcsec, for which they derived a completeness for radio point sources of 85% and a contamination rate of only 3%.
However, at the high angular resolution of FIRST, new problems arise. FIRST is not sensitive to extended radio structures because of a lack of short antennae baselines, and resolves out the extended emission of radio sources. As a result, the total radio luminosity of sources that are larger than a few arcseconds will be systematically low (cf. Becker et al. 1995). In extreme cases, some larger radio sources are missed. These effects introduce systematic biases into the derived radio source sample. In addition, many extended radio sources are split into multiple components by FIRST. Matching routines therefore need to be developed to account for the possible multi-component nature of radio sources.
The first attempt to automate such a routine was by Magliocchetti et al. \[Magliocchetti et al. 1998\], who used a โcollapsing algorithmโ to identify multiโcomponent FIRST sources. They considered all pairs of sources with separations below 3 arcmins, and merged into a single combined source all pairs with separations below $`100\left(S_{\mathrm{tot}}/100\mathrm{m}\mathrm{J}\mathrm{y}\right)^{0.5}`$ arcsec and flux densities within a factor of four of each other. This method is simple and works well for classical doubleโlobed radio sources, but accounts poorly for coreโjet sources or sources with large asymmetries.
Iveziฤ et al. \[Iveziฤ et al. 2002\] improved on this by first cross-correlating all FIRST sources with the SDSS (thereby picking up all sources with a core component) and then adding candidate doubleโlobed radio sources to this sample. These were identified by comparing the mid-points of all FIRST pairs with separations below 90 arcsec with the galaxies in the optical catalogue, and accepting all matches with offsets below 3 arcsec. They estimated that such double sources contribute less than 10% of all radio sources.
McMahon et al. \[McMahon et al. 2002\] carried out a detailed study of the properties of multiโcomponent FIRST sources by comparing isolated pairs of FIRST sources with optical Automated Plate Measuring Machine (APM) scans of the Palomar Observatory Sky Survey (POSS) plates. For coreโjet type sources, where the optical counterpart is associated with one of the radio components, they found that the radio components usually have very different flux densities and that the component with the optical counterpart is usually brighter and is frequently unresolved in the radio. In contrast, if the optical counterpart is located between the two radio components, the two radio components usually have comparable flux densities and similar radio sizes (ie. both are consistent with being radio lobes, not one unresolved core and an extended radio lobe). In this case, the optically identified galaxy is typically located fairly close to the flux-weighted mean position of the two radio components. This information is extremely useful in the identification of multiโcomponent FIRST sources.
Because the main spectroscopic galaxy sample of the SDSS has rather low median redshift ($`z0.1`$), the problems described above associated with identifying extended radio sources will be more severe. This paper thus presents a hybrid method, using information from both NVSS and FIRST in order to take advantage of the strong points of both surveys and avoid the systematic errors that arise in using only one of them. The layout of the paper is as follows. In Section 2, the salient points of the SDSS, NVSS and FIRST surveys are summarised. Section 3 then discusses the cross-matching of these surveys to identify the radio source sample. Section 4 describes how true radio-loud AGN are separated from sources where the radio emission is associated with star formation activity. The local radio luminosity functions of radio-loud AGN and starโforming galaxies are derived in Section 5, and the radio luminosity function of starโforming galaxies is compared to that derived at far-infrared wavelengths. Conclusions are drawn in Section 6. In an accompanying paper \[Best et al. 2005\] the host galaxies of the radioโloud AGN are investigated in detail. Throughout the paper, the values adopted for the cosmological parameters are $`\mathrm{\Omega }_m=0.3`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$, and $`H_0=70`$ km s<sup>-1</sup>Mpc<sup>-1</sup>.
## 2 The radio and optical galaxy samples
### 2.1 The SDSS Spectroscopic Sample
The Sloan Digital Sky Survey (York et al. 2000; Stoughton et al. 2002, and references therein) is an optical imaging (u,g,r,i,z bands) and spectroscopic survey of about a quarter of the extragalactic sky, being carried out at the Apache Point Observatory. The spectroscopic sample considered in this paper is a sample of about 212,000 objects with magnitudes $`14.5<r<17.77`$, spectroscopically confirmed to be galaxies, drawn from the โmain galaxy catalogueโ of the second data release (DR2) of the SDSS. This sample of galaxies is described by Brinchmann et al. \[Brinchmann et al. 2004a\]. The galaxies have a median redshift of $`z0.1`$.
The SDSS spectra cover an observed wavelength range of 3800 to 9200ร
, at an instrumental velocity resolution of about 65km s<sup>-1</sup>. The spectra are obtained through 3-arcsec diameter fibres, which corresponds to 5.7 kpc at a redshift of 0.1; at this redshift the spectra therefore represent a large proportion (up to 50%) of the total galaxy light, while for the very lowest redshift objects they are more dominated by the nuclear emission.
As described in Brinchmann et al. \[Brinchmann et al. 2004a\], a variety of physical parameters for these galaxies have been derived from the photometric and spectroscopic data, and catalogues of these parameters are publically available. These include total stellar masses, mass-to-light ratios, 4000ร
break strengths, H$`\delta `$ absorption measurements and estimates of dust attenuation \[Kauffmann et al. 2003a, Kauffmann et al. 2003b\]; accurate emission line fluxes, after subtraction of the modelled stellar continuum to account for underlying stellar absorption features (Kauffmann et al. 2003c, Tremonti et al. in preparation); galaxy metallicities \[Tremonti et al. 2004\]; parameters measuring optical AGN activity, such as emission line ratios, and galaxy velocity dispersions (hence black hole mass estimates; Kauffmann et al. 2003c, Heckman et al. 2004). These parameters have been adopted for the analyses of this paper: the reader is referred to the papers referenced above for detailed information about the methods used to derive them.
It should be emphasised that the use of the SDSS main galaxy catalogue as the basis sample for this study means that objects classified as โquasarsโ by the automated SDSS classification pipeline (Schlegel et al. in preparation) are excluded. These objects are excluded because of the influence of the direct nonโstellar continuum light from the active nucleus. This affects the observed optical magnitudes, preventing clean magnitudeโlimited samples from being derived, and prohibits the host galaxy parameters discussed above from being accurately determined. The number of AGN excluded in this way is very small ($`\mathrm{}<3`$% out to $`z=0.1`$<sup>1</sup><sup>1</sup>1A search of the SDSS DR2 database reveals only 393 objects classified as โquasarsโ in the redshift range $`0.03<z<0.1`$, compared to 16661 objects classified as emissionโline AGN by Kauffmann et al. \[Kauffmann et al. 2003c\].), since this exclusion only applies to the most luminous Type-I AGN: the SDSS pipeline classifies most low luminosity Type-I AGN as โgalaxiesโ rather than โquasarsโ, and so they are retained. Kauffmann et al. \[Kauffmann et al. 2003c\] estimate that 8% of AGN in the main galaxy sample have broad emission lines, and therefore are strictly Type-Is. They also demonstrate that for these low luminosity objects the non-stellar continuum light has a negligible effect on the physical parameters derived for the host galaxy.
### 2.2 The NVSS and FIRST Radio Surveys
The NVSS \[Condon et al. 1998\] and FIRST \[Becker et al. 1995\] surveys are two radio surveys that have been carried out in recent years using the VLA radio synthesis telescope at a frequency of 1.4 GHz. The NVSS was observed with the array in D-configuration (DnC configuration for the most southerly fields), which provides an angular resolution of 45 arcseconds. This survey covers the entirety of the sky north of $`40^{}`$ declination, down to a limiting point source flux density of about 2.5 mJy. The FIRST observations were carried out in B-array configuration, which provides a much higher angular resolution of $`5`$ arcsec. This survey was designed to study the region of sky that will be observed by the SDSS, and therefore overlaps with this very closely. It reaches a limiting flux density of about 1 mJy for point sources.
## 3 Cross-matching of the SDSS spectroscopic sample using a combination of NVSS and FIRST
The FIRST and NVSS surveys are highly complementary for identifying radio sources associated with nearby galaxies; NVSS provides the sensitivity to largeโscale radio structures required to detect all of the emission from extended radio sources, while FIRST provides the high angular resolution required to reliably identify the host galaxy. In order to identify radio sources associated with galaxies in the SDSS spectroscopic sample, a hybrid method using both radio surveys has been derived. A broad overview of the steps in this process is as follows:
1. SDSS galaxies lying outside the sky coverage of the FIRST survey were excluded. Galaxies close to very bright radio sources were also excluded, because the noise of the NVSS images is much greater in these regions. Finally, galaxies with redshifts below 0.01 were excluded, because at these low redshifts the galaxies are very extended and their optical positions are consequently uncertain.
2. The remaining sample was crossโcorrelated with the NVSS catalogue. A list of candidate galaxies that might be associated with multi-NVSS-component radio sources was derived.
3. These multiโNVSSโcomponent candidates were investigated; by necessity, a small proportion of this analysis had to be done visually rather than through automated procedures. If a galaxy was confirmed to be associated with a multiโNVSSโcomponent source, the integrated flux densities of the NVSS components were summed to provide the radio source flux density.
4. All galaxies matched with a single NVSS source were then crossโcorrelated with the FIRST catalogue. Note, however, that the presence of a FIRST counterpart was not required for a source to be accepted. If there was no FIRST counterpart, then the source was accepted or rejected solely upon its NVSS properties.
5. If a single FIRST counterpart was associated with the NVSS source, then the source was accepted or rejected on the basis of the properties of the FIRST counterpart. For accepted matches, however, the adopted radio flux density was taken from the NVSS data.
6. If multiple FIRST components were associated with the NVSS source, then the source was accepted if it satisfied criteria for a single-component source (with unrelated additional FIRST sources) or for a radio source with multiple FIRST components. Again, the NVSS catalogue was used to provide the most accurate measure of the radio flux density.
The exact criteria for accepting and rejecting matches in the procedures outlined above were tested and refined using MonteโCarlo simulations. Ten catalogues of random sky locations were constructed, over the same sky area as the SDSS survey. Each catalogue contained the same number as positions as the list of SDSS galaxies, and these random catalogues were taken through exactly the same steps of crossโcomparison with the radio data as the SDSS galaxy catalogue. In the subsections that follow, the resulting optimal selection criteria are described, together with the completeness and reliability estimates provided by the MonteโCarlo simulations.
Note that the flux densities adopted for the NVSS sources are true integrated flux densities, rather than the peak flux densities quoted in the NVSS catalogues. The formulae for conversion of peak flux densities to integrated flux densities are provided by Condon et al. \[Condon et al. 1998\]. Only those radio sources with total flux densities (after summing NVSS components if necessary) above 5 mJy are retained. This flux density limit corresponds to approximately 10 times the noise level of the NVSS maps, and is adopted because: (i) at this significance level, all sources should be real and have wellโdetermined positions; (ii) at this flux density limit, the sample is as sensitive to extended singleโcomponent NVSS sources (which will have a lower peak flux density) as it is to point sources, and the sensitivity to multiโcomponent NVSS sources will not be significantly worse (for example, a 5 mJy source composed of two individual components of 2.5 mJy would be found). The 5 mJy limit corresponds to about $`10^{23}`$W Hz<sup>-1</sup> at redshift $`z0.1`$, which is approximately where the radio luminosity function switches from being dominated by star forming galaxies (low luminosities) to being dominated by AGN (high luminosities).
### 3.1 Identification of multi-component NVSS sources
In order to search for possible multi-component NVSS sources, a search was made for multiple sources within a radius of 3 arcmins from each optical galaxy. This distance was selected to be large enough that any genuine multi-component radio source should have at least two matches, but still much smaller than the typical separation of NVSS sources (8-10 arcmins).
Candidate NVSS doubles. For galaxies with two NVSS matches within 3 arcmins, the top panels of Fig 1 compare the offsets of the two NVSS matches from the optical position for SDSS galaxies (left) and for an equivalent number of random positions (right). There are a large number of SDSS galaxies for which the nearer NVSS component lies within 15 arcsec of the optical galaxy; these are predominantly galaxies containing a singleโcomponent NVSS source and the other NVSS source is physically unrelated. Such sources were classified as singleโcomponent matches (see below).
In addition to these, there is a clear excess of SDSS galaxies (compared to random) that have the two NVSS components each offset by 20-50 arcsec from the optical position. For these systems, the flux-weighted mean position of the two NVSS sources is often within 15 arcsec of the optical galaxy (indicated by the diamonds in the upper panel of Fig 1). Candidate NVSS doubles are therefore selected to be sources with both NVSS components closer than 90 arcsec, a flux-weighted mean position closer than 15 arcsec, and the nearer component offset by more than whichever is smaller out of 15 arcsec and the offset of the second source minus 20 arcsec. These selection criteria are indicated by the lines on Fig 1. The 90 arcsec limit is chosen since larger offsets are relatively unlikely and the contamination by random galaxies gets increasingly large beyond this. Even with this limit, there is still significant contamination, but the next step of comparison with FIRST helps to alleviate much of this.
All of these candidate doubles were cross-correlated with the FIRST catalogue. If these are true extended doubles then they may have a central FIRST component associated with a radio source core, and in addition they are likely to be missing flux in the FIRST data due to their extended nature; indeed they may well be undetected by FIRST. If they are not true doubles, but two individual NVSS sources, then it is likely that a single or double FIRST counterpart is present at each NVSS location, with little missing flux. The candidate doubles were thus classified into three categories:
(a) accepted doubles: sources were accepted as NVSS doubles if they either have a FIRST source within 3 arcsec of the optical position, or they satisfy the following three conditions (i) no detected FIRST component (ie. all of the flux is resolved out by FIRST); (ii) both NVSS components lie within 60 arcsec of the SDSS position (larger sources may have additional NVSS components outside of the 3 arcmin limit, and so need to be checked visually); and (iii) the angle NVSS1-SDSS-NVSS2 greater than 135 degrees (ie. consistent with a double radio source with a bend of $`<45`$ degrees).
(b) rejected doubles: those sources with 3 or fewer FIRST components, all further than 15 arcsec from the optical galaxy, and with total flux greater than half of the sum of the two NVSS fluxes, were rejected.
(c) uncertain cases: any sources not satisfying either of the above conditions were classified as uncertain, and referred for visual analysis.
The lower two plots of Figure 1 show the results of this classification of candidate doubles for the SDSS sources and an equivalent number of random positions.
Candidate NVSS triples. Galaxies with 3 NVSS components within 3 arcmins could represent one of four possibilities: (i) a triple radio source associated with the galaxy; (ii) a double radio source associated with the galaxy, together with an unassociated NVSS source; (iii) a single radio source, with two unassociated sources (or an unassociated double source); (iv) 3 unassociated NVSS sources. It is the first two possibilities that are the concern for the multipleโsource analysis.
Comparison between the SDSS galaxies and the random positions (Fig 2) suggests that a source should be classified as a potential triple if all three components are within 120 arcsec, and one of the following three conditions is also satisfied: (i) the flux weighted mean position of all three components is within 15 arcsec of the optical galaxy position; (ii) the flux weighted mean position of the two more distant components is within 15 arcsec of the optical galaxy position \[this for the case where these are the two outer lobes of a radio source, and the nearest component is a feature in the jet of one of the sources\]; (iii) the nearest component is within 15 arcsec of the optical galaxy, with the second and third components both within 90 arcsec and the angle NVSS2-NVSS1-NVSS3 greater than 135 degrees \[this is the case where the nearest component corresponds to the core of the triple: the offset and angle classification requirements distinguish this from a single component source with two unassociated sources\].
Galaxies which satisfied these constraints were investigated using FIRST to accept or reject obvious cases. Galaxies were accepted as triples if they possessed a FIRST source within 3 arcsec of the optical galaxy position. They were rejected if, as for the doubles, they had 3 or fewer FIRST components, all further than 15 arcsec from the optical galaxy, with total flux equal to at least half of the sum of the three NVSS fluxes. The remainder of the galaxies were referred to visual analysis. Figure 2 compares the results of this analysis for the SDSS galaxies and the random sample.
Galaxies that were neither classified as triples nor visually inspected were then investigated to see if they were associated with a double radio source. Each of the three NVSS pairs was checked using the double source analysis described above.
Candidates with 4 or more NVSS matches. Galaxies with 4 or more NVSS sources within 3 arcminutes of the galaxy are likely either to be the host galaxy of a multiple-component radio source, or to lie close to one. No such cases were accepted without visual analysis. Visual analysis was carried out on all galaxies with 5 or more NVSS matches, as well as on galaxies with 4 NVSS matches where either the mean position or the flux-weighted mean position of the 4 NVSS sources was within 30 arcsec of the optical galaxy. All other galaxies were not considered to be quadrupoles, but were examined for potential triples and doubles using the criteria described above.
Overall, as a result of the multiple source analysis, 60 SDSS galaxies were accepted as NVSS multiple-component sources, compared to only 0.3 multiple-component sources predicted from the MonteโCarlo simulations using random positions. This corresponds to less than 1% contamination. A further 277 sources (0.13% of the original sample) required visual analysis, of which 109 were ultimately accepted as genuine sources. This total of 169 accepted multiโNVSSโcomponent sources corresponds to about 6% of the entire SDSS radio source sample.
### 3.2 Single-component NVSS matches
For all galaxies not classified as multiple NVSS sources, Figure 3 compares the distribution of offsets between SDSS galaxies and their nearest NVSS source with the result obtained for random positions. There is a clear excess of sources associated with SDSS galaxies at small radii. Because multiple sources have been removed and the analysis is restricted to brighter sources with well-defined positions, true sources with offsets larger than about 15 arcsec are not expected, but the excess is significant out to at least 100 arcsec. At these large radii the excess is not due to true associations, but rather is the result of the clustering of optical galaxies: on average there are more galaxies within $`100`$ arcsec of an optical galaxy than within the same distance of a random position, and hence there is an increased chance of finding an unassociated radio source. Note that in principle this effect could be accounted for by incorporating an appropriate correlation function into the positions of objects in the random catalogues; however, full knowledge of the environments of radio source hosts would be required to do this properly.
Integrating under the two curves of Figure 3 out to a separation of 15 arcsec gives 2973 matches for the SDSS galaxies and 311 random matches, so the NVSS data alone would suffer from $`10`$% contamination if a 15 arcsec separation were adopted. This falls to $`6`$% contamination at 10 arcsec, but at a cost of reducing the completeness by 10%. The completeness can be improved by including information from the FIRST data. All SDSS galaxies with a single-component NVSS match within 30 arcsec are thus cross-correlated with the FIRST catalogue to determine the number of FIRST components within the 30 arcsec radius.
#### 3.2.1 Sources with no FIRST matches
These sources are either variable radio-loud AGN which have faded between the NVSS and FIRST observations or they are extended radio sources which are resolved out of the FIRST dataset. In either case they should be retained if they are associated with the optical galaxy. These NVSS sources are accepted as matches if they lie within 10 arcsec of the optical position. 134 sources are retained in this way (compared to 8.7 random); this corresponds to approximately 5% of the total radio source sample.
#### 3.2.2 Sources with one FIRST match
For cases with one FIRST source within 30 arcsec, Figure 4 compares the distribution of separations of the FIRST source and the SDSS galaxy with the result obtained for the random catalogue. An excess of FIRST sources with respect to random is seen at all separations; at large separations this is the result of the clustering of optical galaxies, as discussed above. The excess becomes particularly pronounced at separations less than 10 arcsec, but at separations larger than $`3`$ arcsec the contamination of random sources is high. Within 3 arcsec separation, however, the fraction of false identifications is very low, $`\mathrm{}<1`$%. This is lower than has been previously derived for simple SDSSโFIRST comparisons (e.g. Iveziฤ et al. 2002, who found a random contamination of about 9% at 3 arcsec radius). This is because the present analysis is limited to galaxies with NVSS counterparts brighter than 5 mJy.
Given the low contamination rate, all FIRST radio sources with offsets below 3 arcsec can clearly be accepted as matches. If all sources between 3 and 10 arcsec were dropped, however, then the completeness would suffer. Figure 5 compares the offset of the FIRST sources against their projected size along the offset direction, i.e. the product of the deconvolved major axis length and the cosine of the angle between the major axis and the offset vector. A significant fraction of the FIRST sources with offsets between 3 and 10 arcsec are found to be extended sources oriented close to the direction of the offset between the optical and radio position. Only a few of the random positions are associated with radio sources with these properties. The selection procedure was therefore refined to accept those FIRST sources which are either (i) within 3 arcsec, or (ii) offset less than 10 arcsec, oriented within 30 degrees of the offset vector, and offset by less than 75% of the projected major axis length of the source. These selection criteria are illustrated on Figure 5. The addition of the 3โ10 arcsec offset sources significantly reduces the incompleteness of the sample for only a small decrease in the reliability.
#### 3.2.3 Sources with two FIRST matches
For galaxies with two FIRST matches, if the closer of the two matches is within 3 arcsec then the source is accepted under the assumption that it is either a single component source with a nearby unassociated source, or the core of a core-jet source. There are 251 such SDSS galaxies (5.3 random). If neither FIRST source is within 3 arcsec, then it is possible that the two FIRST components are two lobes of the same extended radio source with no core. As discussed earlier, McMahon et al. \[McMahon et al. 2002\] found that in this case the two FIRST sources often have the following properties: (i) they have comparable flux densities; (ii) the flux-weighted mean position of the two sources is close to that of the optical galaxy; (iii) the sizes of the two sources are comparable (ie. both are lobes, not a core and a lobe). None of these conditions on their own is sufficient to classify the source as a double without including a lot of false detections, but the MonteโCarlo simulations show that the combination of the three can be quite powerful. Galaxies with two FIRST matches were accepted as double radio sources if the ratio of the radio source sizes, multiplied by the ratio of the radio source flux densities, multiplied by the offset in arcsec of the fluxโweighted mean position, is less than 5. Figure 6 shows the result of this analysis: 116 double sources are selected among SDSS galaxies, but only 1.2 for random positions, implying that the reliability is about 99%.
It should be noted here that in cases where the galaxy is associated with a single FIRST component, and the other component (or components, for the cases with 3 or more matches discussed below) is genuinely unassociated, the NVSS flux will overestimate the radio luminosity due to the contaminating source. However, only in very rare cases (ie. a faint point source lying nearby a much brighter source) would such a correction be significant and these cases could not reliably be separated from coreโjet type sources without visual analysis.
#### 3.2.4 Sources with three FIRST matches
Galaxies for which there are three FIRST matches within 30 arcsec are accepted if any of the following three conditions are satisfied: (i) the nearest match is within 3 arcsec; (ii) any of the three pairs of sources satisfies the criteria to be accepted as a FIRST double source; (iii) the flux-weighted mean position of all three sources is within 3 arcsec of the optical galaxy position, with the angle subtended by the outer two sources relative to the middle one larger than 135 degrees (ie. the source looks like a straight(ish) triple source). Figure 7 shows the results of this selection.
#### 3.2.5 Sources with four or more FIRST matches
Automated classification of more than three sources cannot be carried out in an efficient and reliable way. For galaxies with four or more matches, the nearest three matches are analysed using the criteria for three-source matches to test whether they may be classified as triples, doubles or singles. All galaxies not accepted in this way are sent for visual analysis (a total of 23 optical galaxies or $`0.01`$% of the SDSS sample).
### 3.3 Repeated matches
The final list of matches was examined to ensure that two different SDSS galaxies were not associated with the same NVSS source. This occurred on 24 occasions and these cases were all examined visually. In two cases, two SDSS galaxies were associated with the same NVSS source and there was no FIRST counterpart. In a further 14 cases, two SDSS galaxies were associated with the same NVSS source which had a single FIRST counterpart which lay close enough to both galaxies. For these 16 objects, the nearer galaxy was accepted as the true match and the other galaxy was removed from the radio source catalogue. There were a further 8 cases where it was found that two galaxies matched the same NVSS source but had distinct FIRST counterparts. In other words, both galaxies were genuine radio sources, but at the lower resolution of NVSS they had been convolved together. In these cases the flux density of the NVSS source was divided between the two galaxies according to the ratio of their integrated FIRST flux densities. If the galaxies still remained above the 5 mJy flux density limit, they were retained in the radio source catalogue.
### 3.4 Completeness and Reliability of the matching procedure
The results of the cross-matching procedure are provided in Table 1. This table gives the number of SDSS galaxies accepted as radio sources compared to the number of cases accepted from the same number of random positions, for each different radio source type. It therefore provides a direct measure of the reliability of each of the criteria defined above. Overall, assuming visual analysis to be 100% reliable, only 30.1 false identifications are expected amongst the final sample of 2712 radio sources. This corresponds to an overall reliability of 98.9%. The most unreliable part of the sample selection is for NVSS sources without a FIRST counterpart. Of these, 6% will be false identifications. This is unavoidable. If sources with no FIRST counterpart are excluded, this reduces the completeness and strongly biases the derived radio source sample by removing 5% of the more extended sources.
The completeness of the sample is more difficult to estimate than the reliability, since the true number of matches expected is unknown. However, various estimates can be made. For galaxies with multiple NVSS components, a comparison of the number of candidate NVSS doubles in the SDSS and random samples with the numbers accepted suggests that the completeness is close to 90%. For the single NVSS component sources, Figure 3 shows that there were 2973 SDSS galaxies with an NVSS source within 15 arcsec, compared to only 311 random galaxies. Assuming that this excess is entirely due to genuine sources and that all true matches lie within 15 arcsec, 2662 genuine singleโcomponent NVSS sources are expected. Table 1 indicates that 2543 single component sources were actually found by the adopted selection procedures, of which about 30 will be false detections. An estimate of the completeness is then (2543โ30) / 2662 $`=`$ 94.4%. Note that this value is conservative because a fraction of the excess matches are likely to be associated with companion galaxies, and so 2662 is an overestimate of the true number of expected matches. Therefore, the overall completeness of the sample likely exceeds 95%.
The values quoted for completeness and reliability are for all types of radio source. There will be a small (but unavoidable) bias against extended sources: the completeness for the singleโcomponent FIRST sources approaches 100%, while that of multiโNVSSโcomponent sources is around 90%. Note that completeness estimates from previous crossโcorrelations with the FIRST catalogue have not taken into account the sources missed because sources with radioโoptical offsets greater than 3 arcsec are excluded ($`3`$% of our final source catalogue) as are sources that are completely resolved out by FIRST (5%). These samples will also miss a fraction of the extended NVSS sources (6%) and the multiโcomponent FIRST sources (6%). These omissions have a severe effect on the completeness of any radio sample derived for the SDSS spectroscopic sample using the FIRST survey alone. The radio luminosities of many sources would also be underestimated using FIRST alone: the distribution of FIRST to NVSS flux density ratios for the final sample of sources is plotted in Figure 8, and shows a long tail to low values. Note, however, that all of these effects are somewhat less important when dealing with the complete imaging catalogue of SDSS, for which the galaxies typically lie at higher redshifts.
### 3.5 The final radio source sample
Details of the final SDSS radio sample of 2712 sources are given in Table 2. This table provides the identification details of each source so that they can be matched against either the original spectra or against the catalogues of derived optical properties released by Brinchmann et al. \[Brinchmann et al. 2004a\]. Also provided are the RA and Dec of each source, the host galaxy redshift, the integrated NVSS flux density and, where there is a central FIRST counterpart, the integrated flux density, radio size and offset from the optical galaxy of the central FIRST component. Each radio source is also given a classification to identify its radio properties. Class 1 sources are singleโcomponent NVSS sources with a single FIRST counterpart. Class 2 sources have a single NVSS match which is resolved into multiple components by FIRST. Class 3 sources have a singleโcomponent NVSS source, but no FIRST counterpart. Class 4 sources have multiple NVSS components. The final column of the table classifies each radio source as a starโforming galaxy or a radioโloud AGN, according to the criteria described in Section 4.
## 4 Definition of the radio-loud AGN sample
The sample of radioโemitting galaxies contains both radioโloud AGN and a population of star forming galaxies. The latter emit at radio wavelengths mostly as a result of the synchrotron emission of particles accelerated in supernova shocks, and their radio luminosity is therefore roughly correlated with their star formation rate: a 1.4 GHz radio luminosity of $`10^{22}`$ W Hz<sup>-1</sup> corresponds to a star formation rate of order $`5M_{}`$ yr<sup>-1</sup> (e.g. Condon 1992 and references therein; Carilli 2001). In order to investigate the host galaxies of these two populations, it is first necessary to separate the radioโloud AGN from the starโforming galaxies.
Starโforming galaxies and AGN are often separated using optical emissionโline properties. Sadler et al. \[Sadler et al. 2002\] used a visual emissionโline classification in their study of radio sources in the 2dFGRS: radioโemitting galaxies without detectable emission lines were classified as radioโloud AGN. Kauffmann et al. \[Kauffmann et al. 2003c\] used the location of a galaxy in the \[OIII\] 5007 / H$`\beta `$ versus \[NII\] 6583 / H$`\alpha `$ emission line diagnostic diagram (Baldwin, Phillips & Terlevich 1981; hereafter BPT) to separate optical AGN from normal star forming galaxies. A key result of the Kauffmann et al. study was that a significant fraction of emission-line selected AGN also have associated star formation. This result means that optical line ratio diagnostics should not be used to identify radioโloud AGN, because star formation activity in galaxies with a radioโquiet active nucleus would give rise to radio emission (and hence a radioโloud classification). In addition, for galaxies which do contain a genuine radioโloud AGN, the radio luminosity associated with the active nucleus will be overestimated if there is a significant contribution of star formation to the radio emission.
Machalski & Condon \[Machalski & Condon 1999\] studied radio galaxies in the LCRS and used farโinfrared to radio flux density ratios and farโinfrared spectral indices to separate the radioโloud AGN and star forming populations. The farโinfrared radio correlation for starโforming galaxies (e.g. Yun, Reddy & Condon 2001) could also be used to correct for the contribution of star formation to the radio luminosities of these systems. This is perhaps the ideal method, but unfortunately the Infrared Astronomical Satellite (IRAS) Faint Source Catalogue is not quite deep enough to allow this to be used for the SDSS galaxies in this paper<sup>2</sup><sup>2</sup>2In fact, the IRAS Faint Source Catalogue is deep enough that the majority of starโforming galaxies with 1.4 GHz radio flux densities of 5 mJy are detected (and lie as expected on the farโinfrared radio correlation), plus a few of the AGN, but the observations are not deep enough that all star forming galaxies are detected. The IRAS data cannot therefore be used as a discriminant between the two subclasses, nor to correct for any star formation contribution to the radio emission of the AGN.. A variety of alternative methods were therefore considered, and a procedure based on the location of a galaxy in the plane of $`D_n(4000)`$ versus $`L_{1.4\mathrm{GHz}}/M_{}`$ was adopted. The $`L_{1.4\mathrm{GHz}}/M_{}`$ ratio provides the radio luminosity per stellar mass of the galaxy and $`D_n(4000)`$ is a fairly accurate indicator of mean stellar age for ages below about a Gyr (at higher ages it is also sensitive to metallicity; cf. Kauffmann et al. 2003b). Thus, star forming galaxies would be expected to occupy a wellโdefined locus in this plane, while radioโloud AGN would be offset to higher radio luminosities. This is demonstrated in the first two panels of Figure 9.
The top panel of Figure 9 shows $`D_n(4000)`$ versus $`L_{1.4\mathrm{GHz}}/M_{}`$ for radioโemitting galaxies that are classified as star forming galaxies using the \[OIII\] 5007 / H$`\beta `$ versus \[NII\] 6583 / H$`\alpha `$ emission line diagnostic diagram. The criteria of Kauffmann et al. \[Kauffmann et al. 2003c\] have been adopted and only galaxies with redshifts in the range $`0.03z0.10`$ have been plotted; the lowest redshifts are excluded because aperture corrections are substantial, whilst beyond $`z=0.1`$ the sensitivity to emission lines is low, hampering classification by emissionโline diagnostics. Radio luminosities have been calculated from the fluxes assuming a radio spectral index of 0.7.
Overlaid on this are theoretical predictions, derived using the Bruzual & Charlot \[Bruzual & Charlot 2003\] stellar synthesis models, for the location of galaxies with different star formation histories. For these models, the radio luminosities have been calculated using the prescription of Hopkins et al. \[Hopkins et al. 2001\]: $`L_{1.4\mathrm{GHz}}=1.8\times 10^{21}(\mathrm{SFR}/M_{})`$W Hz<sup>-1</sup>, where SFR is the average star formation rate in the past $`10^8`$ years. Three of the models are for galaxies with exponentiallyโdecaying star formation rates, of characteristic timescales 1,3 and 5 Gyrs; the tracks indicate how these galaxies move across this plane as they age. A fourth model shows the track for a galaxy with a constant star formation rate. Two further models consider an old (10 Gyr) galaxy which has undergone a recent burst of star formation ($`10^7`$ and $`10^8`$ years ago), as might be the case for a merger-triggered event. Here, the loci of the tracks show what happens if different fractions of the total galaxy mass are converted into stars in the burst. These theoretical tracks largely cover the location of the data points.
The middle left panel of Figure 9 shows the same plot, but now includes all radioโemitting galaxies in this redshift range. The different colours represent different galaxy classifications based upon their locations in the BPT diagram (black โ star forming galaxy; red โ composite systems with both star formation and an AGN; orange โ Seyfert AGN; green โ LINER AGN; purple โ no emission lines). A significant number of the objects classified as AGN based on their emissionโline ratios overlap the starโforming population in this plane. This is interpreted as meaning that these are radioโquiet AGN whose radio emission is due to star formation โ the problem identified above. The dotted line on this plot shows the 3 Gyr exponential star formation model from the top panel. The solid line, 0.225 above this in $`D_n(4000)`$, is the proposed division between radio-loud AGN (above the line, plotted as diamonds) and starโforming galaxies (below the line, plotted as crosses). This cutโoff value was chosen to be most consistent with other methods that could have been adopted for AGNโstarburst separation, as illustrated in the later panels. Using this cut-off, 2215 radio sources are classified as radioโloud AGN, and 497 as starโforming. Note that the plots only show the subset of these at redshifts with $`0.03z0.1`$, to avoid overcrowding the plot and to allow a comparison with emissionโline diagnostic methods; the higher redshift sources fill out more of the plane at larger values of $`L_{1.4\mathrm{GHz}}/M_{}`$, and confirm that the location of the proposed cut at those values is appropriate.
The middle right panel shows the BPT emission line diagnostic diagram. It can be seen that the AGNโstarburst separation defined above (ie. diamonds versus crosses) also makes good sense in this plot: (i) the โcompositeโ galaxies that lie close to the star forming galaxy locus are largely classified as starbursts, whilst those nearer to the AGN locus are predominantly classified as radioโloud AGN; (ii) almost all of the LINERS are classified as radioโloud AGN; (iii) the Seyferts close to the LINER region are mostly classified as radio-loud, whilst those with lower \[NII\] 6583 / H$`\alpha `$ ratios are a mixture of the two classes; (iv) the three star forming galaxies now classified as radio-loud AGN all lie near the boundary with composites.
The lower left panel shows the H$`\alpha `$ / H$`\beta `$ line ratio as a function of galaxy stellar mass. The H$`\alpha `$ / H$`\beta `$ line ratio is an approximate measure of dustโreddening; the dotted line shows the expected value for zero reddening. Star forming galaxies form a tight relation between these parameters, with more massive galaxies being more heavily reddened (cf. Figure 6 of Brinchmann et al. 2004b). Radioโloud AGN deviate from this locus, in the sense of having less reddening (due to less star formation and hence less dust) at a given stellar mass; this diagram indicates that the classification division adopted for $`D_n(4000)`$ versus $`L_{1.4\mathrm{GHz}}/M_{}`$ works well.
The final panel shows the distribution of the galaxies in the $`L_{[\mathrm{OIII}]5007}`$ versus $`L_{1.4\mathrm{GHz}}`$ plane. This relation was considered as a way to separate radioโloud AGN and star-forming galaxies; indeed, it can be seen that for the LINERS and the galaxies without emission lines the division agrees very well with that adopted. However, many of the Seyferts and composites lie on the relation defined by the starโforming galaxies in this plane, but are considerably offset from the star forming locus in all of the other plots. It is for this reason that the final classification was not based upon this relation.
These plots demonstrate the reliability of the AGNโstarburst separation using the $`D_n(4000)`$ versus $`L_{1.4\mathrm{GHz}}/M_{}`$ relation: through comparison of the locations of galaxies on different diagnostics, it is estimated that $`\mathrm{}<1`$% of objects will have been misclassified. The $`D_n(4000)`$ versus $`L_{1.4\mathrm{GHz}}/M_{}`$ relation was also used to estimate and to correct for the star formation contribution to the radio luminosity of galaxies classified as radioโloud AGN: for each of these galaxies the โstar formationโ radio luminosity corresponding its 4000ร
break strength, as estimated by the 3 Gyr exponential star formation track (the dotted line in the middle left panel), was subtracted to obtain a corrected AGN radio luminosity. In no case was this correction larger than 15%.
The radioโloud AGN in the sample exhibit a variety of optical properties; some are classified as optical AGN based upon their emission lines while others are optically inactive. Figure 10 shows the cumulative fractions of the different radio source types as a function of redshift. Out to redshifts $`z0.1`$, the relative numbers of radioโloud AGN with and without emission lines are roughly similar. At higher redshifts the proportion of emission-line AGN decreases rapidly; this is because emission lines such as \[OIII\] 5007 become increasing difficult to detect at higher redshift (only lines brighter than $`10^{5.8}L_{}`$ can be detected at $`z=0.1`$), both because of the increased distance and because the larger physical size of the spectroscopic fibres means that a larger fraction of starlight from the host galaxy is included. This makes it more difficult to pick out the weaker nuclear lines.
## 5 The local radio luminosity function
The local radio luminosity function was derived both for radioโloud AGN and radioโemitting starโforming galaxies out to redshift 0.3. These were calculated in the standard way using the $`1/V_{\mathrm{max}}`$ method \[Schmidt 1968, Condon 1989\], where $`V_{\mathrm{max}}`$ was calculated using the upper and lower redshift limits determined by the joint radio and optical selection criteria, namely a radio cutโoff of 5 mJy and optical cutโoffs of $`14.5<r<17.77`$. An accurate calculation of the exact area of sky within the overlap region of the SDSS DR2 and the FIRST survey is not simple. The absolute normalisation of the derived radio luminosity function has thus been set by normalising the total radio luminosity function to match that derived for the 2dFGRS by Sadler et al. \[Sadler et al. 2002\] over the radio luminosity range $`10^{23}`$ to $`10^{24.5}`$W Hz<sup>-1</sup>, where the errors on the two luminosity function determinations are both small. Note that no correction has been made for incompleteness or misidentification in the radio samples but, as discussed above, it is expected that this will be relatively small.
The radio luminosity functions are tabulated in Table 3. The uncertainties quoted on the luminosity function determination are the statistical Poissonian errors only; these have become so small for the SDSS sample at some luminosities that systematic errors are likely to dominate. One important source of systematic error will be cosmic variance. Another is the separation of AGN and star forming galaxies: for the highest luminosity bin of the star forming galaxies and the lowest luminosity bin of the AGN, this systematic error is likely to be comparable to or larger than the Poissonian uncertainties.
The radio luminosity functions are displayed in Figure 11, along with the results of Sadler et al. \[Sadler et al. 2002\] for the 2dFGRS and Machalski & Godlowski (2000) for the LCRS (both corrected to the cosmology adopted in this paper; note that these determinations are not corrected for incompleteness either). The luminosity function of radioโloud AGN generally agrees well with these previous analyses, although with notably smaller errors. The apparent small mismatches between the Sadler et al. results and those of this paper at $`10^{24.5}`$ and $`10^{25.5}`$W Hz<sup>-1</sup> are likely to be due to cosmic variance.
The luminosity function of star forming galaxies has similar shape to previous measurements, but with slightly higher space densities at high luminosities ($`L_{1.4\mathrm{GHz}}>10^{23}`$ W Hz<sup>-1</sup>). There are three possible differences between the analyses which could account for this. First, the different optical magnitude limits of the different surveys may influence the population of radio starโforming galaxies studied. Second, the combined FIRSTโNVSS radioโoptical crossโcorrelation method adopted here will lead to a more complete sample, particularly for low-redshift star forming galaxies with extended radio emission. Third, the disparity might arise from the contrasting ways in which different analyses treat radioโquiet AGN with associated star formation activity, for which the radio emission is due to the star formation. The technique used in the current paper for separating the star forming and AGN populations would classify such objects as star forming galaxies, but in emission line ratio classifications (as used, for example, by Sadler et al.) they might be classified as AGN. This would lead to previous studies estimating a lower space density of star forming galaxies, particularly at the highest radio luminosities. Note that if star formation dominates the radio and farโinfrared emission of these radioโquiet AGN, then these objects would be expected to lie on the farโinfrared radio correlation for star forming galaxies, despite being classified (by emission line means) as AGN. In this respect it is interesting to consider Figure 10 of Sadler et al. which shows the farโinfrared radio relation for the objects in their sample; many of the objects which lie on the far-infrared radio relation for star forming galaxies, but which have radio luminosities $`L_{1.4\mathrm{GHz}}>10^{23}`$ W Hz<sup>-1</sup>, are indeed classified as AGN. If just some of these objects are truly radio quiet AGN, with radio emission due to associated star formation activity, they could easily account for the small difference between the luminosity function determinations.
It is instructive to compare the radio luminosity function for star forming galaxies with that derived at far-infrared (FIR) wavelengths. Starโforming galaxies show a tight correlation between their radio and farโinfrared luminosities, which Yun et al. \[Yun et al. 2001\] showed for the range of luminosities probed in the current study to be indistinguishable from a linear relation: $`L_{1.4\mathrm{GHz}}/\mathrm{WHz}^1=10^{11.95}L_{60\mu \mathrm{m}}/L_{}`$. The local FIR luminosity function need only be adjusted by this factor to estimate the local radio luminosity function.
Takeuchi et al. \[Takeuchi et al. 2003\] derived the local FIR luminosity function using the IRAS Point Source Catalogue redshift survey (PSCz; Saunders et al. 2000). They fitted the data with an analytic function of the form suggested by Sandage et al. \[Sandage et al. 1979\], namely
$$\varphi (L)=\varphi _{}\left(\frac{L}{L_{}}\right)^{1\alpha }\mathrm{exp}\left(\frac{1}{2\sigma ^2}\left[\mathrm{log}\left(1+\frac{L}{L_{}}\right)\right]^2\right)$$
with $`\alpha =1.23\pm 0.04`$, $`\varphi _{}=(2.60\pm 0.30)\times 10^2h^3`$Mpc<sup>-3</sup>, $`\sigma =0.724\pm 0.01`$ and a characteristic luminosity of $`L_{}=(4.34\pm 0.87)\times 10^8h^2L_{}`$ (where $`h`$ is $`H_0`$ in units of 100 km s<sup>-1</sup>Mpc<sup>-1</sup>). Converting the characteristic luminosity by the factor given above, the equivalent local radio luminosity function is shown as the solid line in Figure 12. This provides a good match to the data at high radio luminosities, supporting the division adopted between AGN and star forming galaxies. The match is less than perfect at lower luminosities, however. The luminosity range covered by the radio observations is not sufficient constrain the parameters for a fit of the above form using the radio data, but the dotted and dashed lines show the effect of doubling the characteristic luminosity or setting the faint end slope $`\alpha `$ to unity (with corresponding changes in $`\varphi _{}`$); these provide a much better fit to the data. However, these represent much larger changes than are allowed by the errors on the fitted FIR parameters or on the radio to FIR conversion.
The difference between the radio and FIR luminosity functions means that the observed radio to FIR correlation stops being a linear relation at low luminosities. This has been suggested before in terms of a steepening of the relation below $`L_{60\mu \mathrm{m}}10^9L_{}`$ (cf. Yun et al. 2001 and references therein) but questions have been raised as to whether the previous studies may be affected by selection biases in the samples under study. Here it is shown that the difference is also present in the luminosity functions. The difference occurs for the lowest luminosity sources, which are generally at the lowest redshifts and therefore may have larger angular sizes, in which case it may be caused by photometric errors due to aperture effects. This is unlikely to be the case, however. The NVSS is sensitive to emission on angular scales out to several arcmins, much larger than potential host galaxies. The IRAS photometry is susceptible to missing extended emission, but for PSCz galaxies this has been corrected for \[Saunders et al. 2000\], and in any case any missing IRAS flux would lower the far-IR to radio ratio, which is in the opposite sense to the observed differences. More likely is that there is a genuine effect at work: either the radio luminosity or the FIR luminosity is not directly proportional to the star formation. One way in which this might occur is if there is an additional lowโlevel contribution to the FIR luminosity of galaxies, for example from dust heated by low mass stars, which is usually swamped by the FIR emission associated with star formation, but becomes significant at low star formation rates (cf. Devereux & Eales 1989 and references therein).
## 6 Conclusions
The main results of this paper are as follows:
* A catalogue of 2712 radio sources has been derived by crossโcorrelating the SDSS spectroscopic sample with a combination of the NVSS and FIRST surveys.
* The use of a hybrid NVSSโFIRST method to identify the radio sources has been highly successful, resulting in a sample with a reliability of 98.9% and a completeness which is estimated to be over 95%.
* The radio sources have been subโdivided into 2215 radioโloud AGN and 497 starโforming galaxies, based upon their location in the plane of 4000ร
break strength versus radio luminosity per unit stellar mass.
* The local radio luminosity functions of radioโloud AGN and star forming galaxies have been derived separately. These are in excellent agreement with previous studies, but with smaller uncertainties.
* The local radio luminosity function of starโforming galaxies has been compared to that derived in the farโinfrared. Differences between the two confirm that the farโIR to radio correlation becomes nonโlinear at low luminosities.
The development of the hybrid NVSSโFIRST method for identification of radio sources represents a large step forward in the study of radio source host galaxies. One note of caution needs to be added: the parameters for acceptance or rejection of sources have been optimised for the SDSS spectroscopic sample, and if the comparison survey has a significantly different sky surface density of objects, the offset parameters for acceptance may need to be modified in order to retain optimal completeness and reliability. The general method, however, can (and where possible, should) be adopted unchanged in future studies.
The sample of radio sources produced will prove invaluable in the study of the host galaxies of radioโloud AGN, due to the large size of the sample and the wealth of information available on the host galaxies from the SDSS. Such analysis is the focus of the accompanying paper.
## Acknowledgements
PNB would like to thank the Royal Society for generous financial support through its University Research Fellowship scheme. The authors thank Jarle Brinchmann, Stephane Charlot, Christy Tremonti and Simon White for making their catalogues available and for useful discussions. The research makes use of the SDSS Archive, funding for the creation and distribution of which was provided by the Alfred P. Sloan Foundation, the Participating Institutions, the National Aeronautics and Space Administration, the National Science Foundation, the U.S. Department of Energy, the Japanese Monbukagakusho, and the Max Planck Society. The research uses the NVSS and FIRST radio surveys, carried out using the National Radio Astronomy Observatory Very Large Array: NRAO is operated by Associated Universities Inc., under co-operative agreement with the National Science Foundation.
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# Slow-roll corrections to inflaton fluctuations on a brane
## I Introduction
Inflation is probably the simplest scenario for the origin of primordial fluctuations in our Universe inflation . Small scale vacuum fluctuations can be stretched to astrophysical scales by an period of accelerated expansion. Inflation provides a test of high-energy physics because the perturbations are generated from very short scales at high energies in the very early universe. These perturbations carry signatures from high energy physics, which can be tested by astronomical observations.
The slow-roll approximation slow-roll is a useful tool to study the fluctuations generated during inflation. If we can neglect the coupling to metric perturbations and the effective mass of the field then the perturbations are described by the fluctuations of a free scalar field in de Sitter spacetime. This gives the familiar result that the power spectrum of scalar field perturbations at horizon-crossing is given by $`(H/2\pi )^2`$. One can then calculate the comoving curvature perturbation which is conserved on super-horizon scales for adiabatic perturbations.
However inflaton perturbations will be coupled to gravity (metric perturbations) at first-order in the slow-roll parameters. In four-dimensional general relativity it is known how to consistently include linear metric perturbations by working in terms of the gauge invariant combination of scalar field and curvature perturbations, the so-called Mukhanov-Sasaki variable, which obeys a simple wave equation Mukhanov . Gravitational effects are negligible at small scales and high energies, where perturbations can be normalised to the usual Bunch-Davies vacuum state. On large scales (super-horizon scales) the comoving curvature perturbation is conserved allowing one to relate observations of temperature anisotropies in the cosmic microwave background to high energy vacuum fluctuations during inflation. Exact solutions are known for the special case of power-law inflation in general relativity which generalise the de Sitter result and have been used to calculate first-order slow-roll corrections in more general inflation models S-L .
In this paper, we develop a new way to derive slow-roll corrections based on a slow-roll expansion about de Sitter spacetime. In four-dimensional general relativity we show how to recover the usual first-order slow-roll corrections. Our method may be useful when one cannot derive an exact solution and the background spacetime is given as a perturbation about de Sitter spacetime.
We then apply our method to inflation in the brane world model. New ideas in the string theory suggest that our observable universe is a 4-dimensional hypersurface, or brane, in a higher dimensional bulk spacetime Review . The simplest example of this model is the Randall-Sundrum model where there is a brane embedded in a 5-dimensional anti-de Sitter (AdS) spacetime RS99 . An AdS spacetime has a characteristic curvature scale $`\mu `$ associated with the negative cosmological constant in the bulk. The spacetime shrinks exponentially away from the brane and this geometry effectively compactifies the 5-dimensional spacetime with the effective size $`\mu ^1`$. On large length scales $`L>\mu ^1`$, 4-dimensional Einstein gravity is recovered, while on small scales, the gravity becomes 5-dimensional GT . In the early universe when the Hubble horizon is smaller than $`\mu ^1`$, we expect significant effects from higher dimensional bulk spacetime. Indeed, the Friedmann equation is modified from the conventional 4-dimensional theory for $`H\mu 1`$ BDEL . This modification of Friedmann can provide a novel model for inflation MWBH ; TMM ; binf .
In Ref.MWBH , the amplitude of the curvature perturbation is calculated by taking into account the modification of the Friedmann equation. This work has been extended to include higher order corrections in slow-roll parameters Liddle and the formula has been widely used to confront this model with the observations obs . But to derive these formulae for the spectrum of the primordial curvature perturbations the effect of coupling to five-dimensional gravity has been neglected and in particular it is assumed that the power spectrum of inflaton perturbations at horizon crossing is given by $`(H/2\pi )^2`$. This assumption is only valid to zeroth order in slow-roll parameters. At first order the inflaton perturbations will be coupled to metric perturbations. In the brane world, metric perturbations live in the 5-dimensional spacetime, and thus we must check if 5-dimensional effects change the result of conventional 4-dimensional theory. Especially, at small scales/high energies, the 5-dimensional effects could be large.
The first attempt to study the backreaction due to metric perturbations was made in Ref KLMW . There, perturbations are solved perturbatively in slow-roll parameters. We should emphasize that this is the only possible way to perform the calculations analytically. If the background spacetime of the brane deviates from de Sitter spacetime, we cannot solve the bulk metric perturbations analytically. In contrast to four-dimensional general relativity, there are no other exact solutions known for the perturbation equations. Thus we must develop a new approach to calculate the effect of slow-roll corrections. In this paper, we extend earlier studies and investigate the backreaction due to higher-dimensional perturbations using a slow-roll expansion.
The structure of the rest of the paper is as follows. In section II, we describe our new approach to derive first order slow-roll corrections in a conventional 4-dimensional cosmology. In section III, we review an inflation model in the Randall-Sundrum brane-world driven by an inflaton field on the brane. In section IV, we derive the equations that govern the coupled system of inflaton fluctuations on the brane and metric perturbations in the bulk. In section V, the first order corrections to the inflaton fluctuations on the brane are solved. In section VI, we discuss the implications of our result for the brane world inflation model.
## II Slow-roll expansion of scalar perturbations in 4D cosmology
### II.1 Background spacetime
We consider an inflaton $`\varphi `$ whose potential energy density $`V(\varphi )`$ drives inflation. In the conventional 4-dimensional general relativity described by the metric
$$ds^2=dt^2+a(t)^2\delta _{ij}dx^idx^j,$$
(1)
the Friedmann equation and the equation of motion for the homogeneous field, $`\varphi `$, are given by
$`H^2`$ $`=`$ $`{\displaystyle \frac{\kappa _4^2}{3}}\left({\displaystyle \frac{1}{2}}\dot{\varphi }^2+V(\varphi )\right),`$ (2)
$`\ddot{\varphi }`$ $`+`$ $`3H\dot{\varphi }={\displaystyle \frac{dV}{d\varphi }},`$ (3)
where $`H=\dot{a}/a`$, $`\kappa _4=8\pi G_4`$ and $`G_4`$ is the 4D gravitational coupling constant. A dot indicates a derivative with respect to cosmic time, $`t`$. Slow-roll parameters are defined by
$$ฯต\frac{\dot{H}}{H^2},\eta \frac{\ddot{\varphi }}{H\dot{\varphi }}.$$
(4)
Slow-roll inflation is described by small values of $`ฯต`$ and $`\eta `$.
### II.2 Slow-roll corrections to inflaton fluctuations
The inhomogeneous inflaton fluctuation, $`\delta \varphi `$, is coupled to the metric perturbations. In the Longitudinal gauge, the perturbed metric is written as
$$ds^2=(1+2\mathrm{\Psi })dt^2+a(t)^2(1+2\mathrm{\Phi })\delta _{ij}dx^idx^j.$$
(5)
The coupled equations for $`\delta \varphi `$, $`\mathrm{\Psi }`$ and $`\mathrm{\Phi }`$ can be simplified by using Mukhanov-Sasaki variable defined by Mukhanov
$$u=a\left(\delta \varphi \frac{\dot{\varphi }}{H}\mathrm{\Psi }\right).$$
(6)
Expanding $`u`$ by Fourier modes, the wave equation for $`u`$ is given by
$`{\displaystyle \frac{d^2u_k}{d\tau ^2}}+\left(k^2{\displaystyle \frac{1}{z}}{\displaystyle \frac{d^2z}{d\tau ^2}}\right)u_k=0,`$ (7)
where $`z(a\dot{\varphi })/H`$ and $`\tau `$ is a conformal time defined as
$$\tau =\frac{dt}{a(t)}.$$
(8)
In the case of the slow-roll inflation, the mass term in Mukhanov-Sasaki equation (7) can be approximated as
$`{\displaystyle \frac{1}{z}}{\displaystyle \frac{d^2z}{d\tau ^2}}={\displaystyle \frac{1}{\tau ^2}}\left(2+6ฯต3\eta +๐ช(\eta ^2,ฯต^2)\right),`$ (9)
up to first order of the slow-roll parameters. Then Eq. (7) can be expressed as
$`{\displaystyle \frac{d^2u_k}{d\tau ^2}}+\left(k^2{\displaystyle \frac{1}{\tau ^2}}(2+6ฯต3\eta )\right)u_k=0.`$ (10)
Usually, the appropriately normalized solution with the correct asymptotic behavior at small scales is obtained by solving Eq. (10) directly as
$`u_k(\tau )={\displaystyle \frac{\sqrt{\pi }}{2}}e^{i(\nu +1/2)\pi /2}(\tau )^{1/2}H_\nu ^{(1)}(k\tau ),`$ (11)
where $`\nu =3/2+2ฯต\eta `$ and $`H_\nu ^{(1)}`$ is the Hankel function of the first kind of order $`\nu `$. Here we assumed the Bunch-Davies vacuum state where perturbations stay in Minkowski vacuum at small scales. Equation (11) is an exact solution of the perturbation equation (10) only if the slow-roll parameters $`ฯต`$ and $`\eta `$ are constant. However their variation in a Hubble time is second-order and hence of higher-order in the slow-roll expansion. Thus we can take $`ฯต`$ and $`\eta `$ to be evaluated around the time of horizon-crossing.
We are interested in the asymptotic form of the solution well outside the horizon. Taking the limit $`k\tau 0`$ yields the asymptotic form of $`u_k`$;
$`u_ke^{i(\nu 1/2)\pi /2}2^{\nu 3/2}{\displaystyle \frac{\mathrm{\Gamma }(\nu )}{\mathrm{\Gamma }(3/2)}}{\displaystyle \frac{1}{\sqrt{2k}}}(k\tau )^{\nu +1/2}.`$ (12)
Expanding the gamma function in Eq. (12), we get
$`u_ke^{i(\nu 1/2)\pi /2}\left\{1+(2ฯต\eta )(2\gamma \mathrm{ln}2)\right\}{\displaystyle \frac{1}{\sqrt{2k}}}(k\tau )^{12ฯต+\eta },`$ (13)
where we have used the formula for the poly-Gamma function
$`\psi (3/2){\displaystyle \frac{\mathrm{\Gamma }^{}(3/2)}{\mathrm{\Gamma }(3/2)}}=2\gamma 2\mathrm{ln}2,`$ (14)
where $`\gamma `$ is an Euler number.
The quantity that is related to observables today is the the power spectrum of the curvature perturbation given by
$`๐ซ_{}^{1/2}(k)=\sqrt{{\displaystyle \frac{k^3}{2\pi ^2}}}\left|{\displaystyle \frac{u_k}{z}}\right|.`$ (15)
From Eqs. (9) and (10), it can be shown that, at large scale, the time dependences of $`u_k`$ and $`z`$ are the same, that is, $`๐ซ_{}^{1/2}`$ is constant. Note that the constancy of $``$ in the large scale limit does not depend on the slow-roll approximation, but holds for any adiabatic perturbation. Thus this comoving curvature perturbation can be related to the perturbation in the radiation density on large scales long after inflation has ended.
For the model with a monotonous potential, the following relation holds:
$`|z|={\displaystyle \frac{a|\dot{\varphi }|}{H}}={\displaystyle \frac{2}{\kappa _4^2}}{\displaystyle \frac{a}{H}}\left|{\displaystyle \frac{dH}{d\varphi }}\right|,`$ (16)
and conformal time can be evaluated up to the first order in slow-roll parameters as
$`\tau ={\displaystyle \frac{1}{aH}}(1+ฯต).`$ (17)
Thus the power spectrum of the curvature perturbation is given by
$`๐ซ_{}^{1/2}`$ $`=`$ $`[1(2C+1)ฯต+C\eta ]{\displaystyle \frac{\kappa _4^2}{4\pi }}\left\{{\displaystyle \frac{H^2}{|dH/d\varphi |}}\right\}_{k=aH},`$ (18)
where $`C=2+\mathrm{ln}2+\gamma 0.73`$. The terms proportional to the slow-roll parameters are called Stewart-Lyth correction S-L .
### II.3 Perturbing about de Sitter spacetime
In this subsection, we reproduce the usual slow-roll corrections in a perturbative approach which does not require any exact solution of the perturbation equation other than that in a de Sitter spacetime. This will be more suited to extension to the case of brane-world gravity.
At zeroth order in slow-roll parameters, the spacetime is described by the de Sitter spacetime. Thus we can expand the spacetime from de Sitter spacetime. The scale factor is expanded as
$$a(t)=a^{(0)}(t)+a^{(1)}(t)+๐ช(ฯต^2),a^{(0)}(t)=\mathrm{exp}(Ht).$$
(19)
Accordingly, the Mukhanov-Sasaki variable is expanded as
$`u_k(\tau )=u_k^{(0)}(\tau )+u_k^{(1)}(\tau )+๐ช(ฯต^2),`$ (20)
where $`u_k^{(0)}a\delta \varphi ^{(0)}`$ and $`u_k^{(1)}a(\delta \varphi ^{(1)}(\dot{\varphi }/H)\mathrm{\Psi })`$. Substituting Eq. (20) into Eq. (10), the zeroth order equation is given by
$`{\displaystyle \frac{d^2u_k^{(0)}}{d\tau ^2}}+\left(k^2{\displaystyle \frac{2}{\tau ^2}}\right)u_k^{(0)}=0.`$ (21)
Since we expect that the effects of the deviation from de-Sitter spacetime are insignificant at small scales, the form of $`u_k^{(0)}`$ is determined by demanding a Bunch-Davies vacuum
$`u_k^{(0)}(\tau )=A(\tau )^{1/2}H_{3/2}^{(1)}(k\tau ),`$ (22)
where $`A=(\sqrt{\pi }/2)e^{i\theta }`$ and the phase $`\theta `$ is fixed so that $`u_k(\tau )(1/\sqrt{2k})e^{ik\tau }`$.
Next, we must solve $`u_k^{(1)}`$ sourced by this zeroth order solution
$`{\displaystyle \frac{d^2u_k^{(1)}}{d\tau ^2}}+\left(k^2{\displaystyle \frac{2}{\tau ^2}}\right)u_k^{(1)}{\displaystyle \frac{1}{\tau ^2}}(6ฯต3\eta )u_k^{(0)}=0.`$ (23)
If we impose the boundary conditions (i) $`u_k^{(1)}(\tau )`$ is negligible in the limit $`\tau \mathrm{}`$, and (ii) $`u_k^{(1)}(\tau )`$ does not diverge faster than $`u_k^{(0)}(\tau )`$ in the limit $`\tau 0`$, then we find that the solution is given by,
$`u_k^{(1)}`$ $`=`$ $`C_1(k\tau )^{1/2}J_{3/2}(k\tau )+C_2(k\tau )^{1/2}H_{3/2}^{(1)}(k\tau ),`$
$`C_1`$ $`=`$ $`{\displaystyle \frac{\pi i}{2}}(6ฯต3\eta )A{\displaystyle _{\mathrm{}}^\tau }๐\tau ^{}{\displaystyle \frac{1}{\tau ^{}}}\left\{H_{3/2}^{(1)}(k\tau ^{})\right\}^2,`$
$`C_2`$ $`=`$ $`{\displaystyle \frac{\pi i}{2}}(6ฯต3\eta )A{\displaystyle _{\mathrm{}}^\tau }๐\tau ^{}{\displaystyle \frac{1}{\tau ^{}}}H_{3/2}^{(1)}(k\tau ^{})J_{3/2}(k\tau ^{}),`$ (24)
where $`J_\nu `$ is the Bessel function of the order $`\nu `$.
We take the limit $`k\tau 0`$ and compare the asymptotic form with Eq. (13). Using the small arguments limit of the Bessel functions
$`J_{3/2}(x)\left({\displaystyle \frac{x}{2}}\right)^{3/2}{\displaystyle \frac{1}{\mathrm{\Gamma }(5/2)}},H_{3/2}^{(1)}(x)i{\displaystyle \frac{\mathrm{\Gamma }(3/2)}{\pi }}\left({\displaystyle \frac{2}{x}}\right)^{3/2},`$ (25)
we can show that the zeroth order Mukhanov-Sasaki variable approaches to
$`u_k^{(0)}(\tau )4i{\displaystyle \frac{\mathrm{\Gamma }(\frac{3}{2})}{\pi }}A(2k)^{1/2}(k\tau )^1.`$ (26)
Next, we must evaluate the asymptotic form of the first order Mukhanov-Sasaki variable. Using Eq. (25), $`C_1`$ is evaluated as
$`C_1`$ $``$ $`4i(2ฯต\eta )A{\displaystyle \frac{\mathrm{\Gamma }(3/2)^2}{\pi }}{\displaystyle \frac{1}{(k\tau )^3}}.`$ (27)
We should be careful in evaluating the asymptotic behavior of $`C_2`$ because sub-leading terms are comparable to the contribution from $`C_1`$. Using
$`J_{3/2}(x)=\sqrt{{\displaystyle \frac{2}{\pi x}}}\left({\displaystyle \frac{\mathrm{sin}x}{x}}\mathrm{cos}x\right),J_{3/2}(x)=\sqrt{{\displaystyle \frac{2}{\pi x}}}\left(\mathrm{sin}x+{\displaystyle \frac{\mathrm{cos}x}{x}}\right),`$ (28)
the integral for $`C_2`$ in Eq. (24) can be evaluated as
$`{\displaystyle _{\mathrm{}}^\tau }๐\tau ^{}{\displaystyle \frac{1}{\tau ^{}}}H_{3/2}^{(1)}(k\tau ^{})J_{3/2}(k\tau ^{}){\displaystyle \frac{2i}{3\pi }}\text{Ci}(2k\tau )+{\displaystyle \frac{14i}{9\pi }},`$ (29)
where Ci is the integrated cosine function defined as
$`Ci(x){\displaystyle _x^{\mathrm{}}}{\displaystyle \frac{\mathrm{cos}t}{t}}๐t.`$ (30)
For small $`k\tau `$, the integrated cosine function can be expressed as
$`Ci(2k\tau )\gamma +\mathrm{ln}2+\mathrm{ln}(k\tau ).`$ (31)
Therefore, the asymptotic form of $`C_2`$ is given by
$`C_2`$ $``$ $`(2ฯต\eta )A\left(\gamma +\mathrm{ln}2+\mathrm{ln}(k\tau )+{\displaystyle \frac{7}{3}}\right).`$ (32)
Then we obtain the asymptotic form of $`u_k`$ for $`k\tau 0`$ up to the first order in slow-roll parameters
$`u_k(\tau )`$ $``$ $`(i)e^{i\theta }\left\{1+(2ฯต\eta )(2\gamma \mathrm{ln}2)\right\}\left\{1(2ฯต\eta )\mathrm{ln}(k\tau )\right\}{\displaystyle \frac{1}{\sqrt{2k}}}(k\tau )^1,`$ (33)
where we have used the fact that $`A=(\sqrt{\pi }/2)e^{i\theta }`$. This should be compared with Eq. (13). There appears a logarithmic term which diverges for $`k\tau 0`$. However, if we can renormalize this divergence by rewriting the logarithmic term as
$`1(2ฯต\eta )\mathrm{ln}(k\tau )(k\tau )^{2ฯต+\eta }.`$ (34)
we see that Eq. (33) is consistent with Eq. (13).
Indeed, the logarithmic divergence for $`k\tau 0`$ in $`u_k`$ does not show up in the spectrum of the curvature perturbation. In order to see this, we expand the curvature perturbation as
$`๐ซ_{}^{1/2}=\{๐ซ_{}^{1/2}\}^{(0)}+\{๐ซ_{}^{1/2}\}^{(1)}+๐ช(ฯต^2).`$ (35)
On the other hand, by the definition of the curvature perturbation (15), we can write the spectrum of curvature perturbation up to the first order in slow-roll parameters as
$`๐ซ_{}^{1/2}\sqrt{{\displaystyle \frac{k^3}{2\pi }}}\left|{\displaystyle \frac{u_k^{(0)}}{z^{(0)}}}+{\displaystyle \frac{u_k^{(0)}}{z^{(0)}}}\left({\displaystyle \frac{u_k^{(1)}}{u_k^{(0)}}}{\displaystyle \frac{z^{(1)}}{z^{(0)}}}\right)\right|,`$ (36)
where we also expanded $`z(a\dot{\varphi })/H`$ as
$`z=z^{(0)}+z^{(1)}+๐ช(ฯต^2).`$ (37)
Since there is a difficulty to define the curvature perturbation in de Sitter spacetime, we concentrate on the ratio between the zeroth order and the first order of the curvature perturbation. By comparing Eq. (35) to (36), the ratio is given by
$`{\displaystyle \frac{\{๐ซ_{}^{1/2}\}^{(1)}}{\{๐ซ_{}^{1/2}\}^{(0)}}}={\displaystyle \frac{u_k^{(1)}}{u_k^{(0)}}}{\displaystyle \frac{z^{(1)}}{z^{(0)}}}.`$ (38)
In order to evaluate Eq. (38), we must obtain $`z^{(0)}`$ and $`z^{(1)}`$, that is, we must solve Eq. (9) perturbatively. Substituting Eq. (37) into Eq. (9), the equation for $`z`$ at zeroth order is given by
$`{\displaystyle \frac{d^2z^{(0)}}{d\tau ^2}}={\displaystyle \frac{2}{\tau ^2}}z^{(0)}.`$ (39)
If we consider only the growing mode, the zeroth order solution for $`z^{(1)}`$ can be obtained as
$`z^{(0)}=B\tau ^1,`$ (40)
where $`B`$ is an integration constant. This zeroth order solution gives a source term in the equation for $`z`$ at first order;
$`{\displaystyle \frac{d^2z^{(1)}}{d\tau ^2}}={\displaystyle \frac{2}{\tau ^2}}z^{(1)}+{\displaystyle \frac{(6ฯต3\eta )}{\tau ^2}}z^{(0)}.`$ (41)
The growing mode solution for the first order $`z^{(1)}`$ is given by
$`z^{(1)}=(2ฯต\eta )B\tau ^1\mathrm{ln}(k\tau )+BD\tau ^1,`$ (42)
where $`D`$ is another integration constant. Then we get
$`{\displaystyle \frac{z^{(1)}}{z^{(0)}}}=(2ฯต\eta )\mathrm{ln}(k\tau )+D.`$ (43)
This logarithmic divergence term exactly cancels the logarithmic divergence term in $`u_k`$;
$`{\displaystyle \frac{u_k^{(1)}}{u_k^{(0)}}}=(2\gamma \mathrm{ln}2\mathrm{ln}(k\tau ))(2ฯต\eta ).`$ (44)
From Eqs. (43) and (44) we obtain
$`{\displaystyle \frac{\{๐ซ_{}^{1/2}\}^{(1)}}{\{๐ซ_{}^{1/2}\}^{(0)}}}=C(2ฯต\eta )D,`$ (45)
where $`C=2+\mathrm{ln}2+\gamma 0.73`$ is again a numerical constant. We cannot determine $`D`$ in this approach, which comes from the difficulty to define curvature perturbation in pure de Sitter spacetime. However, we can still fix $`D`$ as follows. Neglecting the logarithmic term, which is canceled by the contribution from $`u_k`$, the solution for $`z`$ is written as
$$z=B(1+D)\tau ^1.$$
(46)
This must be compared with the definition of $`z`$
$$z=\frac{a|\dot{\varphi }|}{H}\frac{|\dot{\varphi }|}{H^2}(1+ฯต)\tau ^1,$$
(47)
where the solution for $`a`$ up to the first order was used. Then we can identify $`B=|\dot{\varphi }|/H^2`$ and $`D=ฯต`$. Then Eq. (45) agrees with the Stewart-Lyth correction given by Eq. (18).
## III Slow-roll inflation in Randall-Sundrum brane world
In this section, we apply our perturbative approach to the brane-world model. We consider the simplest version of brane-world inflation model based on the Randall-Sundrum model. We will consider a single brane embedded in a 5-dimensional AdS spacetime. We assume that the inflaton $`\varphi `$ is confined to the brane while gravity can propagate in the whole 5-dimensional spacetime MWBH .
The 5-dimensional metric describing this model is given by BDEL
$$ds^2=dy^2N(y,t)^2dt^2+A(y,t)^2\delta _{ij}dx^idx^j,$$
(48)
where
$`A(y,t)`$ $`=`$ $`a(t)\left[\mathrm{cosh}\mu y\left(1+{\displaystyle \frac{\kappa ^2\rho }{6\mu }}\right)\mathrm{sinh}\mu y\right],`$
$`N(y,t)`$ $`=`$ $`\mathrm{cosh}\mu y\left(1{\displaystyle \frac{\kappa ^2\rho }{6\mu }}(2+3w)\right)\mathrm{sinh}\mu y,`$ (49)
$$\rho =\frac{1}{2}\dot{\varphi }^2+V(\varphi ),P=\frac{1}{2}\dot{\varphi }^2V(\varphi ),$$
(50)
and $`w=P/\rho `$. The brane is located at $`y=0`$ and the inflaton is confined to this hypersurface. On the brane, the Friedmann equation and the equation of motion for the scalar field are given by
$`H^2`$ $`=`$ $`{\displaystyle \frac{\kappa _4^2}{3}}\rho +{\displaystyle \frac{\kappa ^4}{36}}\rho ^2,`$ (51)
$`\ddot{\varphi }`$ $`+`$ $`3H\dot{\varphi }={\displaystyle \frac{dV}{d\varphi }},`$ (52)
where $`\kappa _4^2=\kappa ^2\mu `$, $`\kappa ^2=8\pi G_5`$ and $`G_5`$ is 5D gravitational coupling. We can define slow-roll parameters in the same way as the conventional cosmology, Eq. (4).
Unfortunately, the background metric (48) is not in general a separable function with respect to $`y`$ and $`t`$. Thus we cannot solve the metric perturbations analytically. In order to solve for the $`y`$-dependence of the bulk gravitons and to study the time-dependence of the perturbations on the brane, we will expand about the special case of a de Sitter spacetime on the brane. This corresponds to the background solution to zeroth order in a slow-roll expansion. For a de Sitter brane, AdS bulk gives a separable form for the bulk metric Kaloper :
$$ds^2=dy^2+N^2(y)\left[dt^2+a^2(t)\delta _{ij}dx^idx^j\right],$$
(53)
where
$`a(t)`$ $`=`$ $`e^{Ht},`$ (54)
$`N(y)`$ $`=`$ $`{\displaystyle \frac{H}{\mu }}\mathrm{sinh}\mu (y_\mathrm{h}|y|),.`$ (55)
and $`y=\pm y_\mathrm{h}`$ are Cauchy horizons Kaloper , with
$$y_\mathrm{h}=\frac{1}{\mu }\mathrm{coth}^1\left(\sqrt{1+\left(\frac{H}{\mu }\right)^2}\right).$$
(56)
It is often useful to work in terms of the conformal bulk-coordinate $`z=๐y/N(y)`$:
$$z=\mathrm{sgn}(y)H_o^1\mathrm{ln}\left[\mathrm{coth}\frac{1}{2}\mu (y_\mathrm{h}|y|)\right].$$
(57)
The Cauchy horizon is now at $`|z|=\mathrm{}`$, and the brane is located at $`z=\pm z_\mathrm{b}`$, with
$$z_\mathrm{b}=\frac{1}{H}\mathrm{sinh}^1\frac{H}{\mu }.$$
(58)
The line element, Eq. (53), becomes
$$ds^2=N^2(z)\left[dt^2+dz^2+\mathrm{e}^{2Ht}d\stackrel{}{x}^2\right],$$
(59)
where
$$N(z)=\frac{H}{\mu \mathrm{sinh}(H|z|)}.$$
(60)
## IV Equations for bulk metric perturbations and inflaton perturbations on the brane
In this section, we derive the basic equations for the coupled Mukhanov-Sasaki variable on the brane and bulk metric perturbations following Ref.KLMW .
### IV.1 Master variable for perturbations in AdS bulk
In the background spacetime given by Eq. (59) bulk metric perturbations can be solved using the master variable Mukohyama ; Kodama . The perturbed metric is given by
$$ds^2=N(z)^2\left[(1+2A_{yy})dz^2+2A_ydtdz(1+2A)dt^2+a^2(1+2)\delta _{ij}dx^idx^j\right].$$
(61)
In the special case of a de Sitter brane in the AdS bulk, the metric variables are written by the master variable $`\mathrm{\Omega }`$ as
$`A`$ $`=`$ $`{\displaystyle \frac{a^1N^3}{6}}\left(2\mathrm{\Omega }^{\prime \prime }3{\displaystyle \frac{N^{}}{N}}\mathrm{\Omega }^{}+\ddot{\mathrm{\Omega }}\mu ^2N^2\mathrm{\Omega }\right),`$ (62)
$`A_y`$ $`=`$ $`a^1N^3\left(\dot{\mathrm{\Omega }}^{}{\displaystyle \frac{N^{}}{N}}\dot{\mathrm{\Omega }}\right),`$ (63)
$`A_{yy}`$ $`=`$ $`{\displaystyle \frac{a^1N^3}{6}}\left(\mathrm{\Omega }^{\prime \prime }3{\displaystyle \frac{N^{}}{N}}\mathrm{\Omega }^{}+2\ddot{\mathrm{\Omega }}+\mu ^2N^2\mathrm{\Omega }\right),`$ (64)
$`R`$ $`=`$ $`{\displaystyle \frac{a^1N^3}{6}}\left(\mathrm{\Omega }^{\prime \prime }\ddot{\mathrm{\Omega }}2\mu ^2N^2\mathrm{\Omega }\right).`$ (65)
From the perturbed 5-dimensional Einstein equation, we can derive the equation for $`\mathrm{\Omega }`$
$`\ddot{\mathrm{\Omega }}3H\dot{\mathrm{\Omega }}\left(\mathrm{\Omega }^{\prime \prime }3{\displaystyle \frac{N^{}}{N}}\mathrm{\Omega }^{}\right)+{\displaystyle \frac{k^2}{a^2}}\mathrm{\Omega }\mu ^2N^2\mathrm{\Omega }=0.`$ (66)
Solutions of the master equation can be separated into eigenmodes of the time-dependent equation on the brane and bulk mode equation:
$$\mathrm{\Omega }(t,y;\stackrel{}{x})=d^3\stackrel{}{k}๐m\alpha _m(t)u_m(z)e^{i\stackrel{}{k}.\stackrel{}{x}},$$
where
$`\ddot{\alpha }_m3H\dot{\alpha }_m+\left[m^2+{\displaystyle \frac{k^2}{a^2}}\right]\alpha _m`$ $`=`$ $`0,`$ (67)
$`u_m^{\prime \prime }3{\displaystyle \frac{N^{}}{N}}u_m^{}+\mu ^2N^2u_m+m^2u_m`$ $`=`$ $`0.`$ (68)
Note that the Hubble damping term $`3H\dot{\alpha }_m`$ has the โwrong signโ, i.e., this is not the standard wave equation for a scalar field in four-dimensions.
If we write $`\alpha _m=a^2\phi _m`$ and work in terms of the conformal time $`\tau =1/(aH)`$, the time-dependent part of the wave equation (67) can be rewritten as
$$\frac{d^2\phi _m}{d\tau ^2}+\left[k^2\frac{2(m^2/H^2)}{\tau ^2}\right]\phi _m=0.$$
This is the same form of the time-dependent mode equation commonly given for a massive scalar field in de Sitter spacetime. The general solution is given by
$$\phi _m(\eta ;\stackrel{}{k})=\sqrt{k\tau }B_\nu (k\tau ),\nu ^2=\frac{9}{4}\frac{m^2}{H^2},$$
(69)
where $`B_\nu `$ is a linear combination of Bessel functions of order $`\nu `$. The solutions oscillate at early-times/small-scales for all $`m`$, with an approximately constant amplitude while they remain within the de Sitter event horizon ($`kaH`$). โHeavy modesโ, with $`m>\frac{3}{2}H`$, continue to oscillate as they are stretched to super-horizon scales, but their amplitude rapidly decays away, $`|u_m^2|a^3`$. But for โlight modesโ with $`m<\frac{3}{2}H`$, the perturbations become over-damped at late-times/large-scales ($`kaH`$), and decay more slowly: $`|u_m^2|a^{2\nu 3}`$.
### IV.2 Mukhanov-Sasaki equation on the brane
Now we introduce a scalar field fluctuation on the brane. We expand the scalar field perturbation in terms of slow-roll parameters;
$$\delta \varphi =\delta \varphi ^{(0)}+\delta \varphi ^{(1)}+\mathrm{}$$
(70)
The 0-th order of the scalar field fluctuation obeys the following equation of motion,
$$\delta \ddot{\varphi }^{(0)}+3H\delta \dot{\varphi }^{(0)}+\frac{k^2}{a^2}\delta \varphi ^{(0)}=0.$$
(71)
The metric perturbations are generated by the 0-th order fluctuation of the scalar field through the induced Einstein equations on the brane Deffayet ,
$`3H\dot{\mathrm{\Psi }}3H^2\mathrm{\Phi }+{\displaystyle \frac{k^2}{a^2}}\mathrm{\Psi }`$ $`=`$ $`{\displaystyle \frac{\kappa _{4,\mathrm{eff}}^2}{2}}(\dot{\varphi }\dot{\delta \varphi }_0+V^{}\delta \varphi ^{(0)})+{\displaystyle \frac{\kappa _4^2}{2}}\delta \rho _E,`$ (72)
$`H\mathrm{\Phi }\dot{\mathrm{\Psi }}`$ $`=`$ $`{\displaystyle \frac{\kappa _{4,\mathrm{eff}}^2}{2}}\dot{\varphi }\delta \varphi ^{(0)}{\displaystyle \frac{\kappa _4^2}{2}}\delta q_E,`$ (73)
$`\ddot{\mathrm{\Psi }}3H\dot{\mathrm{\Psi }}+H\dot{\mathrm{\Phi }}+3H^2\mathrm{\Phi }{\displaystyle \frac{1}{3}}{\displaystyle \frac{k^2}{a^2}}(\mathrm{\Psi }+\mathrm{\Phi })`$ $`=`$ $`{\displaystyle \frac{\kappa _{4,\mathrm{eff}}^2}{2}}(\dot{\varphi }\dot{\delta \varphi }_0V^{}\delta \varphi ^{(0)})+{\displaystyle \frac{\kappa _4^2}{6}}\delta \rho _E,`$ (74)
$`a^2(\mathrm{\Psi }+\mathrm{\Phi })`$ $`=`$ $`\kappa _4^2\delta \pi _E,`$ (75)
where
$$A(y=0,t)=\mathrm{\Phi }(t),R(y=0,t)=\mathrm{\Psi }(t)$$
(76)
$`\kappa _4^2\delta \rho _E`$ $`=`$ $`{\displaystyle \frac{k^4a^5}{3}}\mathrm{\Omega }`$ (77)
$`\kappa _4^2\delta q_E`$ $`=`$ $`{\displaystyle \frac{k^2a^3}{3}}\left(\dot{\mathrm{\Omega }}H\mathrm{\Omega }\right),`$ (78)
$`\kappa _4^2\delta \pi _E`$ $`=`$ $`{\displaystyle \frac{a^3}{2}}\left(\ddot{\mathrm{\Omega }}H\mathrm{\Omega }+{\displaystyle \frac{k^2a^2}{3}}\mathrm{\Omega }\right),`$ (79)
and
$$\kappa _{4,\mathrm{eff}}=\kappa _4\frac{N^{}}{N}|_{y=0}.$$
(80)
The contributions $`\delta \rho _E,\delta q_E`$ and $`\delta \pi _E`$ come from the projected 5D Weyl tensor and these describe the effect of the bulk gravitational perturbations SMS . The metric fluctuations in turn affect the dynamics of the first order scalar field perturbation
$$\ddot{\delta \varphi }^{(1)}+3H\dot{\delta \varphi }^{(1)}+\frac{k^2}{a^2}\delta \varphi ^{(1)}=V^{\prime \prime }\delta \varphi ^{(0)}3\dot{\varphi }\dot{\mathrm{\Psi }}+\dot{\varphi }\dot{\mathrm{\Phi }}2V^{}\mathrm{\Phi }.$$
(81)
In order to evaluate the effect from metric perturbations, it is useful to use Mukhanov-Sasaki variable $`Q`$ as in the conventional cosmology;
$$Q=\delta \varphi \frac{\dot{\varphi }}{H}\mathrm{\Psi }.$$
(82)
In terms of slow-roll expansion, we have $`Q^{(0)}=\delta \varphi ^{(0)}`$ and $`Q^{(1)}=\delta \varphi ^{(1)}(\dot{\varphi }/H)\mathrm{\Psi }`$. Then using the induced Einstein equations, Eqs.(72), (74) and (75), we can derive the equation for $`Q^{(1)}`$;
$$\ddot{Q}^{(1)}+3H\dot{Q}^{(1)}+\frac{k^2}{a^2}Q^{(1)}=V^{\prime \prime }Q^{(0)}6\dot{H}Q^{(0)}+J,$$
(83)
where
$`J`$ $`=`$ $`{\displaystyle \frac{\kappa _4^2\dot{\varphi }}{3H}}\left(k^2\delta \pi _E+\delta \rho _E\right)`$ (84)
$`=`$ $`{\displaystyle \frac{\dot{\varphi }}{H}}{\displaystyle \frac{k^2a^3}{6}}\left(\ddot{\mathrm{\Omega }}H\dot{\mathrm{\Omega }}+{\displaystyle \frac{k^2}{a^2}}\mathrm{\Omega }\right).`$
The equation is the same as the standard 4-dimensional cosmology except for the term $`J`$, which describes the corrections from the 5-dimensional bulk perturbations. Because $`J`$ contains the 5-dimensional quantity $`\mathrm{\Omega }`$ we must solve the bulk equation for $`\mathrm{\Omega }`$ to evaluate the effects.
### IV.3 Boundary condition for $`\mathrm{\Omega }`$
In order to solve $`\mathrm{\Omega }`$, we must specify the boundary condition for $`\mathrm{\Omega }`$. we rewrite the expressions of $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$, Eq.(62) and (65), as HK
$`\mathrm{\Psi }`$ $`=`$ $`{\displaystyle \frac{a^1N^3}{6}}\left[3{\displaystyle \frac{N^{}}{N}}3H(\dot{\mathrm{\Omega }}H\mathrm{\Omega })a^2\mathrm{\Delta }\mathrm{\Omega }\right],`$ (85)
$`\mathrm{\Phi }`$ $`=`$ $`{\displaystyle \frac{a^1N^3}{6}}\left[3{\displaystyle \frac{N^{}}{N}}3\ddot{\mathrm{\Omega }}+6H\dot{\mathrm{\Omega }}3H^2\mathrm{\Omega }+2a^2\mathrm{\Delta }\mathrm{\Omega }\right].`$ (86)
where
$$=\mathrm{\Omega }^{}\frac{N^{}}{N}\mathrm{\Omega }.$$
(87)
Substituting these expressions into the induced Einstein equations (72)-(75), we obtain the equations written only by $``$ and $`\delta \varphi ^{(0)}`$:
$`3H\dot{}k^2a^2`$ $`=`$ $`\kappa ^2a(\dot{\varphi }\dot{\delta \varphi ^{(0)}}+V^{}(\varphi )\delta \varphi ^{(0)}),`$ (88)
$`\dot{}`$ $`=`$ $`\kappa ^2a\dot{\varphi }\delta \varphi ^{(0)},`$ (89)
$`\ddot{}+2H\dot{}`$ $`=`$ $`\kappa ^2a(\dot{\varphi }\dot{\delta \varphi ^{(0)}}V^{}(\varphi )\delta \varphi ^{(0)}).`$ (90)
These equations can be thought as the boundary conditions for $`\mathrm{\Omega }`$. Combining the junction conditions, Eqs.(88)-(90), we get an evolution equation for $``$;
$$\ddot{}\left(H+2\frac{\ddot{\varphi }}{\dot{\varphi }}\right)\dot{}+k^2a^2=0.$$
(91)
This is consistent with the equation for scalar field equation Eq.(71).
## V Perturbative solutions
We must solve the coupled equations Eqs. (66) for $`\mathrm{\Omega }`$ and Eq. (83) for $`Q`$. Introducing dimensionless quantities
$$Q(t)=Ha(t)^1u(\tau ),\mathrm{\Omega }(z,t)=\kappa ^2\dot{\varphi }H^1\omega (z,\tau ),$$
(92)
the coupled equations are written as
$`k^2\tau ^2\left(\ddot{\omega }+{\displaystyle \frac{4}{\tau }}\dot{\omega }+k^2\omega \right)`$ $`=`$ $`\omega ^{\prime \prime }+3{\displaystyle \frac{\mathrm{cosh}Hz}{\mathrm{sinh}Hz}}\omega ^{}+{\displaystyle \frac{1}{\mathrm{sinh}^2Hz}}\omega ,`$ (93)
$`\dot{}_\omega `$ $`=`$ $`aH^2u,_\omega =\left(\omega ^{}+{\displaystyle \frac{\mathrm{cosh}Hz}{\mathrm{sinh}Hz}}\omega \right)_{z=z_b},`$ (94)
$`\ddot{u}+k^2u{\displaystyle \frac{1}{\tau ^2}}(2+6ฯต3\eta )u`$ $`=`$ $`J_u,J_u=\beta ^2k^2\tau ^2\left(\ddot{\omega }+{\displaystyle \frac{2}{\tau }}\dot{\omega }+k^2\omega \right),`$ (95)
where a dot denotes a derivative with respect to $`\tau `$ and
$$\beta ^2=\frac{\kappa ^2\dot{\varphi }^2}{6H}.$$
(96)
At the leading order in slow-roll parameters, $`\beta ^2`$ can be written as
$$\beta ^2=\frac{1}{3}ฯต\frac{H}{\mu }\left(1+\left(\frac{H}{\mu }\right)^2\right)^{1/2}.$$
(97)
Thus $`\beta ^2`$ is essentially the slow-rolling parameter and it controls the strength of coupling between inflaton perturbation and gravitational perturbations in the bulk. We solve the coupled equations perturbatively in terms of small $`\beta ^2`$.
### V.1 Zeroth order solutions
At the zeroth order where $`\beta ^2=0`$, the solution for $`u`$ is given by
$$u^{(0)}=C_1(k\tau )^{1/2}J_{3/2}(k\tau )+C_2(k\tau )^{1/2}J_{3/2}(k\tau ).$$
(98)
Then $`_\omega `$ becomes
$`(\tau )=C_1H\sqrt{{\displaystyle \frac{2}{\pi }}}{\displaystyle \frac{\mathrm{cos}(k\tau )}{k\tau }}+C_2H\sqrt{{\displaystyle \frac{2}{\pi }}}{\displaystyle \frac{\mathrm{sin}(k\tau )}{k\tau }}.`$ (99)
This gives the boundary condition for $`\omega `$. The solution for $`\omega `$ in the bulk subject to this condition is obtained as KLMW
$`\omega ^{(0)}(z,\tau )`$ $`=`$ $`2C_1{\displaystyle \underset{\mathrm{}=0}{\overset{\mathrm{}}{}}}(1)^{\mathrm{}}\left(2\mathrm{}+{\displaystyle \frac{1}{2}}\right){\displaystyle \frac{(\mathrm{sinh}Hz_b)Q_2\mathrm{}(\mathrm{cosh}Hz)}{\mathrm{sinh}HzQ_2\mathrm{}^1(\mathrm{cosh}Hz_b)}}(k\tau )^{3/2}J_{2\mathrm{}+1/2}(k\tau ),`$ (100)
$`+2C_2{\displaystyle \underset{\mathrm{}=0}{\overset{\mathrm{}}{}}}(1)^{\mathrm{}}\left(2\mathrm{}+{\displaystyle \frac{3}{2}}\right){\displaystyle \frac{\mathrm{sinh}Hz_bQ_{2\mathrm{}+1}(\mathrm{cosh}Hz)}{\mathrm{sinh}HzQ_{2\mathrm{}+1}^1(\mathrm{cosh}Hz_b)}}(k\tau )^{3/2}J_{2\mathrm{}+3/2}(k\tau ),`$
where the identities
$`\mathrm{cos}(x)`$ $`=`$ $`\sqrt{2\pi }{\displaystyle \underset{\mathrm{}=0}{\overset{\mathrm{}}{}}}(1)^{\mathrm{}}\left(2\mathrm{}+{\displaystyle \frac{1}{2}}\right)x^{\frac{1}{2}}J_{2\mathrm{}+1/2}(x),`$ (101)
$`\mathrm{sin}(x)`$ $`=`$ $`\sqrt{2\pi }{\displaystyle \underset{\mathrm{}=0}{\overset{\mathrm{}}{}}}(1)^{\mathrm{}}\left(2\mathrm{}+{\displaystyle \frac{3}{2}}\right)x^{\frac{1}{2}}J_{2\mathrm{}+3/2}(x),`$ (102)
were used.
At large scales $`k\tau 0`$, the dominant contribution comes from $`\mathrm{}=0`$ mode. On the other hand, on small scales $`k\tau \mathrm{}`$, all modes becomes comparable and we need to take into account an infinite ladder of the modes. This means that gravity becomes 5-dimensional at small scales.
In practice, we must approximate the infinite sum to proceed the calculations. We first check the identity Eqs. (101) and (102) to see if we can approximate the infinite summation by introducing a cut-off $`\mathrm{}_c`$ into the summation. From Fig. 1, we can see that if we increase the cut-off $`\mathrm{}_c`$, the identity is satisfied for large $`k\eta `$, i.e. on small scales. This implies that as long as we start from a finite time $`k\tau _i`$, we can approximate the infinite ladder of the modes by introducing sufficiently large $`\mathrm{}_c`$.
Fig. 2 shows the bulk solution for $`\omega (z,t)`$ with introducing sufficiently large cut-off $`\mathrm{}_c`$. The solution is localized near the brane and decays towards the horizon $`z\mathrm{}`$. This is a bound state that is supported by an oscillation of the inflaton fluctuation on the brane. This kind of bound state generally appears in coupled brane and bulk oscillators couple . A toy example is shown in Appendix. A key point here is that, in this case, the bound state is a summation of many different eigenstates of different eigenvalues (Eq. (100)). This fact becomes crucial in the analysis of the next order solution.
### V.2 First order solutions
Now it is possible to calculate the next order equation for $`u^{(1)}`$
$`{\displaystyle \frac{d^2u_k}{d\tau ^2}}+\left(k^2{\displaystyle \frac{1}{\tau ^2}}(2+6ฯต3\eta )\right)u_k=J_u,`$ (103)
where $`J_u`$ describes the effect of the back reaction from the bulk perturbations. We can use the 0-th order solution to evaluate $`J_u`$ as
$`J_u`$ $`=`$ $`{\displaystyle \frac{2}{3}}ฯตk^2C_1{\displaystyle \underset{\mathrm{}=0}{\overset{\mathrm{}}{}}}(1)^{\mathrm{}}\left(2\mathrm{}+{\displaystyle \frac{1}{2}}\right)(2\mathrm{};H\mu )\left(2\mathrm{}(2\mathrm{}1)(k\tau )^{\frac{3}{2}}J_{2\mathrm{}+\frac{1}{2}}+2(k\tau )^{\frac{1}{2}}J_{2\mathrm{}+\frac{3}{2}}(k\tau )\right)`$ (104)
$``$ $`{\displaystyle \frac{2}{3}}ฯตk^2C_2{\displaystyle \underset{\mathrm{}=0}{\overset{\mathrm{}}{}}}(1)^{\mathrm{}}\left(2\mathrm{}+{\displaystyle \frac{3}{2}}\right)(2\mathrm{}+1;H\mu )\left(2\mathrm{}(2\mathrm{}+1)(k\tau )^{\frac{3}{2}}J_{2\mathrm{}+\frac{3}{2}}+2(k\tau )^{\frac{1}{2}}J_{2\mathrm{}+\frac{5}{2}}(k\tau )\right),`$ (105)
where
$`(n;H\mu )={\displaystyle \frac{H}{\mu }}\left(1+\left({\displaystyle \frac{H}{\mu }}\right)^2\right)^{1/2}{\displaystyle \frac{Q_n(\mathrm{cosh}Hz_b)}{Q_n^1(\mathrm{cosh}Hz_b)}}.`$ (106)
The quantity $`(n;H\mu )`$ controls the amplitude of corrections to Mukhanov-Sasaki equations from the bulk over the change of the energy scales of the inflation.
In order to evaluate $`J_u`$, we need to introduce a cut-off in the summation at sufficiently large $`\mathrm{}`$. Fig. 3 shows the behaviour of $`J_u`$ against the change of the cut-off $`\mathrm{}_c`$. A good feature here is that the behavior of $`J_u`$ for small $`k\tau `$ does not change even if we increase the cut-off. Thus we can reproduce the correct behaviour of $`J_u`$ by a finite summation of modes as long as we are considering a finite time interval.
#### V.2.1 Large scales
On large scales $`k\tau 0`$, $`\mathrm{}=0`$ mode in $`C_1`$ mode dominates, which corresponds to a $`m^2=2H^2`$ mode. Thus we can approximate the infinite ladder of the modes by a single mode on super horizon scales. This indicates that, at large scales, gravity looks four-dimensional. Then we can easily show that $`J_u`$ is suppressed for $`k\tau 0`$ and the Mukhanov-Sasaki equation becomes completely the same as the conventional cosmology. Thus we can show the conservation of the curvature perturbation $``$ on large scales in the same way as conventional cosmology LMSW .
#### V.2.2 Small scales
At low energies $`H/\mu 1`$, $`(n;H/\mu )`$ can be approximated as
$$(n;H/\mu )=\left(\frac{H}{\mu }\right)^2\left(\gamma +\psi (n+1)+\mathrm{log}(H/\mu )\mathrm{log}2\right),$$
(107)
where we assumed $`n`$ is not large. Thus, the source terms is well suppressed by the term $`(n;H/\mu )`$ at low energies. However, at sufficiently small scales, large $`\mathrm{}`$ modes become important and the approximation (107) does not hold. Then we could still get an effect on very sub-horizon scales ($`k\mu ^1H`$). In this case, we need to introduce a large cut-off in the summation of $`\mathrm{}`$ and it is technically difficult to perform a calculation.
At high energies, the amplitude of $`(n;H\mu )`$ becomes large as $`H\mu `$ becomes large, but, at sufficient high energies $`H\mu \mathrm{}`$, $`(n;H\mu )`$ becomes independent of $`H\mu `$ as seen from Fig.4. Indeed, we can obtain the asymptotic form of $`(l;H/\mu )`$ for $`H/\mu \mathrm{}`$ as
$$(l;H\mu )\frac{1}{n+1}.$$
(108)
In the following, we consider this limit. In this high energy limit, $`J_u`$ is well fitted as
$`J_u{\displaystyle \frac{2ฯต}{3}}k^2A\left[C_1(k\tau )^{1/2}\mathrm{cos}(k\tau +\phi )C_2(k\tau )^{1/2}\mathrm{sin}(k\tau +\phi )\right],`$ (109)
between $`140<k\tau <40`$ where $`A=0.4`$ and $`\phi =0.9`$. Then, the equation of motion for the first order Mukhanov variable is given as
$`{\displaystyle \frac{d^2u_k^{(1)}}{d\tau ^2}}`$ $`+`$ $`\left(k^2{\displaystyle \frac{2}{\tau ^2}}\right)u_k^{(1)}{\displaystyle \frac{1}{\tau ^2}}(6ฯต3\eta )u_k^{(0)}J_u(\tau )=0.`$ (110)
By using the asymptotic behavior of the Bessel function at small scale (large $`k\tau `$), the third term behaves like $`(k\tau )^2\mathrm{sin}(k\tau )`$, while the forth term behaves as $`(k\tau )^{1/2}\mathrm{sin}(k\tau )`$. Therefore, at least at small scales, the effect from the bulk metric perturbations dominates the effect from the standard corrections to the de Sitter geometry. Thus we will neglect the third term. The general solutions are given by the linear combination of $`(k\tau )^{1/2}J_{3/2}(k\tau )`$ and $`(k\tau )^{1/2}J_{3/2}(k\tau )`$. By choosing the initial conditions so that $`u_k(\tau _i)=u_k^{(0)}(\tau _i)`$ at $`\tau =\tau _i`$, we find the following form of the solution,
$`u_k^{(1)}=D_1(k\tau )^{1/2}J_{\frac{3}{2}}(k\tau )+D_2(k\tau )^{1/2}J_{\frac{3}{2}}(k\tau ),`$ (111)
where $`D_1`$ and $`D_2`$ are given by
$`D_1`$ $`=`$ $`{\displaystyle \frac{\pi }{2}}{\displaystyle _{k\tau _i}^{k\tau }}d(k\tau ^{})(k\tau ^{})^{\frac{1}{2}}J_{\frac{3}{2}}(k\tau ^{})J_u(\tau ^{}),`$
$`D_2`$ $`=`$ $`{\displaystyle \frac{\pi }{2}}{\displaystyle _{k\tau _i}^{k\tau }}d(k\tau ^{})(k\tau ^{})^{\frac{1}{2}}J_{\frac{3}{2}}(k\tau ^{})J_u(\tau ^{}).`$ (112)
For specifying the behavior of the first order Mukhanov variable, we must evaluate $`D_1`$ and $`D_2`$. Using the asymptotic form for Bessel functions at small scale, $`D_1`$ and $`D_2`$ are well approximated as
$`D_1`$ $``$ $`{\displaystyle \frac{2ฯต}{3}}A\sqrt{{\displaystyle \frac{\pi }{2}}}{\displaystyle _{k\tau _i}^{k\tau }}d(k\tau ^{})(k\tau ^{})^{\frac{1}{2}}\mathrm{sin}(k\tau ^{})\left[C_1\mathrm{cos}(k\tau ^{}+\phi )C_2\mathrm{sin}(k\tau ^{}+\phi )\right],`$
$`D_2`$ $``$ $`{\displaystyle \frac{2ฯต}{3}}A\sqrt{{\displaystyle \frac{\pi }{2}}}{\displaystyle _{k\tau _i}^{k\tau }}d(k\tau ^{})(k\tau ^{})^{\frac{1}{2}}\mathrm{cos}(k\tau ^{})\left[C_1\mathrm{cos}(k\tau ^{}+\phi )C_2\mathrm{sin}(k\tau ^{}+\phi )\right].`$ (113)
Then, on small scales, the first order solution is given by
$`u_k^{(1)}((F(\tau )F(\tau _i))\mathrm{cos}(k\tau )+(G(\tau )G(\tau _i))\mathrm{sin}(k\tau ),`$ (114)
where
$`F(\tau )`$ $`=`$ $`{\displaystyle \frac{ฯตAC_1\sqrt{\pi }}{3}}\left[S\left({\displaystyle \frac{2\sqrt{k\tau }}{\sqrt{\pi }}}\right)\mathrm{cos}\phi \left(C\left({\displaystyle \frac{2\sqrt{k\tau }}{\sqrt{\pi }}}\right){\displaystyle \frac{2}{\sqrt{\pi }}}\sqrt{k\tau }\right)\mathrm{sin}\phi \right]`$ (115)
$`+{\displaystyle \frac{ฯตAC_2\sqrt{\pi }}{3}}\left[S\left({\displaystyle \frac{2\sqrt{k\tau }}{\sqrt{\pi }}}\right)\mathrm{sin}\phi \left(C\left({\displaystyle \frac{2\sqrt{k\tau }}{\sqrt{\pi }}}\right){\displaystyle \frac{2}{\sqrt{\pi }}}\sqrt{k\tau }\right)\mathrm{cos}\phi \right],`$
$`G(\tau )`$ $`=`$ $`{\displaystyle \frac{ฯตAC_1\sqrt{\pi }}{3}}\left[\left(C\left({\displaystyle \frac{2\sqrt{k\tau }}{\sqrt{\pi }}}\right)+{\displaystyle \frac{2}{\sqrt{\pi }}}\sqrt{k\tau }\right)\mathrm{cos}\phi S\left({\displaystyle \frac{2\sqrt{k\tau }}{\sqrt{\pi }}}\right)\mathrm{sin}\phi \right]`$ (116)
$`+{\displaystyle \frac{ฯตAC_2\sqrt{\pi }}{3}}\left[\left(C\left({\displaystyle \frac{2\sqrt{k\tau }}{\sqrt{\pi }}}\right)+{\displaystyle \frac{2}{\sqrt{\pi }}}\sqrt{k\tau }\right)\mathrm{sin}\phi S\left({\displaystyle \frac{2\sqrt{k\tau }}{\sqrt{\pi }}}\right)\mathrm{cos}\phi \right],`$
where $`S`$ and $`C`$ are Fresnel functions.
We see that the first order perturbations grows like $`\sqrt{k\tau }\sqrt{k\tau _i}`$. Then if we formally take the limit $`k\tau _i\mathrm{}`$, the first order corrections diverge. Thus our perturbative approach breaks down. The amplitude of the zeroth order oscillation of inflaton fluctuations are significantly affected by the first order corrections.
We should take care in interpreting this result for the amplitude. In a toy model of a coupled boundary and bulk oscillators described in Appendix A, this change of amplitude due to the first order perturbations is merely caused by the breakdown of the perturbative expansion. In the toy model, the coupling to the bulk oscillator just changes the phase of the brane oscillator. In that case we can renormalize the first-order perturbation so that the first-order corrections appear only in the phase of the oscillations and do not have a large effect on the amplitude. However, in the case of inflaton fluctuations, we cannot do this kind of renormalization. This is due to the phase $`\phi `$ in the source term of the first order equation (see Appendix A). The phase originates from the fact that the zeroth order oscillation cannot be matched by a single bulk eigenmode with the same frequency as the brane oscillator and we need an infinite ladder of modes. Thus we can say that the effects on the amplitude from first order corrections are not artificial effects of our perturbative approach.
In conventional cosmology, the amplitude of inflaton oscillations $`u`$ remains constant, so we can impose initial conditions on any scale far inside the horizon. However, in the brane world case, the coupling to the bulk metric perturbations changes the amplitude of the zeroth order inflaton oscillation $`u`$, so the effect crucially depends on the initial conditions. In general, classically, we can also impose arbitrary initial conditions for $`\mathrm{\Omega }`$. Indeed, it is always possible to add homogeneous solutions which satisfy the boundary condition given by
$$=0.$$
(117)
Then we find an infinite tower of massive modes starting from $`m^2=9H^2/4`$. Arbitrary initial conditions for $`\mathrm{\Omega }`$ can be satisfied by an appropriate summation of these massive modes. These massive modes also affect the evolution of inflaton fluctuations $`u`$ HK .
We have tried to solve the coupled equations for inflaton fluctuations and master variable directly using a numerical method Hiramatsu . If we begin with the initial condition for $`\omega `$ given by Eq.(100), the numerical solution for $`u`$ well agrees with our perturbative solutions as long as perturbations remain valid. We have also tried using different initial conditions for $`\mathrm{\Omega }`$ and find that the effects on the amplitude of $`u`$ depend on the initial conditions for $`\omega `$ in the bulk.
The initial conditions for $`u`$ and $`\omega `$ must be determined by quantum theory on small scales. Thus we must quantise the coupled system of the inflaton fluctuations $`u`$ and the master variable $`\omega `$ consistently. This is in contrast to the conventional cosmology where we can specify the vacuum for $`u`$ by neglecting the gravitational effects far inside the horizon. This means that the assumption that the power spectrum of inflaton perturbations at horizon crossing is given by $`(H/2\pi )^2`$ could be invalid and we may have significant effects on the amplitude of perturbations from the backreaction due to the bulk metric perturbations.
## VI Conclusion
In this paper we have studied the effect of metric perturbations upon inflaton fluctuations during inflation, at first-order in slow-roll parameters $`ฯต`$ and $`\eta `$, which describe the dimensionless slope and curvature of the potential. If we neglect the slope and curvature of the inflaton potential then we obtain the familiar results for free field fluctuations in de Sitter spacetime, with a scale invariant power spectrum on large (super-horizon) scales. We take this as our zeroth-order result in a slow-roll expansion.
In four-dimensional general relativity we were able to calculate corrections to the field evolution perturbatively to first-order in a slow-roll expansion, including linear metric perturbations. As far as we are aware this is the first time the slow-roll corrections have been calculated in the manner. We reproduce the familiar slow-roll corrections usually derived from Lyth and Stewartโs exact solution to the linear perturbation equations in power-law inflation.
On a four-dimensional brane-world, embedded in a five-dimensional bulk, there are no exact solutions for cosmological perturbations (for a vacuum bulk described by Einstein gravity) except for the case of an exactly de Sitter brane. Thus the only way to calculate slow-roll corrections is perturbatively in a slow-roll expansion. We have calculated the leading order bulk metric perturbations sourced by the zeroth-order inflaton fluctuations on the brane. We find that inflaton fluctuations support an infinite tower of discrete bulk perturbations, with negative effective mass-squared.
Including the effect of the metric perturbations as an inhomogeneous source term in the wave equation for the first-order inflaton fluctuations we find that the effect of bulk metric perturbations becomes small on large scales, and we recover the usual result that the comoving curvature perturbation becomes constant outside the horizon.
However at small scales (or early times for a given comoving wavelength) the effect of bulk metric perturbations cannot be neglected. We are able to give an approximate solution for inflaton fluctuations at high energies and on sub-horizon scales using a truncated tower of bulk modes. This shows that the bulk metric perturbations change the amplitude of inflaton field fluctuations on the brane. By including a large number of bulk modes we can model this effect for many oscillations, but ultimately this change of amplitude becomes a large effect leading to a breakdown of our perturbative analysis.
It is not surprising in some ways that we see a large effect at small scales as these are high momentum modes which are expected to be strongly coupled to the bulk. Nonetheless this invalidates the usual assumption that gravitational effects are small far inside the cosmological horizon. It seems necessary to consistently solve for the coupled evolution of brane and bulk modes. We numerically tried to solve this problem and verified the validity of our perturbative approach as long as perturbations remain good. But it was also found that the change of the amplitude depends on the initial conditions for bulk metric perturbations. Detailed analysis of numerical solutions go beyond the scope of the present paper and they will be presented in a separate paper Hiramatsu . In order to give definite predictions for the amplitude of scalar perturbations in high energy inflation, we must specify the quantum vacuum state for coupled inflaton fluctuations and metric perturbations consistently and determine initial conditions. For this purpose, it would be useful to study the quantum theory of the toy model for a coupled bulk-brane oscillators in more details where we can consistently quantise a coupled system couple .
Our result implies the possibility that the assumption that the power spectrum of inflaton perturbations at horizon crossing on a brane is given by $`(H/2\pi )^2`$ could be invalid and we may have significant effects on the amplitude of perturbations from the backreaction due to the bulk metric perturbations.
###### Acknowledgements.
SM is grateful to the ICG, Portsmouth for their hospitality when this work was initiated and DW is grateful to the Maeda Lab, Waseda for their hospitality during its continuation. KK is grateful to T. Hiramatsu for the numerical analysis. KK is supported by the Particle Physics and Astronomy Research Council. SM is supported by the Grant-in-Aid for Scientific Research Fund (Young Scientists (B) 17740154). This work was also supported by PPARC grant PPA/V/S/2001/00544.
## Appendix A Toy model for coupled bulk-brane system
In this appendix, we present a simple toy model for a coupled brane and bulk oscillators. Let us consider a toy model for a brane field $`q(t)`$ and a bulk field $`\varphi `$ in Minkowski bulk, which satisfy
$`\ddot{q}+\mu ^2q`$ $`=`$ $`\beta \varphi ,`$
$`\ddot{\varphi }`$ $`=`$ $`\varphi ^{\prime \prime }m^2\varphi ,\varphi ^{}(y=0)={\displaystyle \frac{\beta }{2}}q.`$ (118)
We solve the equations perturbatively in terms of small $`\beta `$. Without coupling, the zeroth order solution for $`q`$ is given by
$$q^{(0)}(t)=C_1\mathrm{cos}(\mu t)+C_2\mathrm{sin}(\mu t).$$
(119)
If we assume $`m>\mu `$, the 0-th order solution for $`\varphi `$ is obtained as
$$\varphi ^{(0)}=\frac{\beta }{2\sqrt{m^2\mu ^2}}\left(C_1\mathrm{cos}(\mu t)+C_2\mathrm{sin}(\mu t)\right)e^{\sqrt{m^2\mu ^2}y}.$$
(120)
Note that the bulk field has a negative effective mass-squared and decays towards $`y\mathrm{}`$. This is a normalizable bound state supported by an oscillation of $`q(t)`$ on the brane. The equation for the next order $`q^{(1)}(t)`$ is given by
$$\ddot{q}^{(1)}=\mu ^2q^{(1)}+\frac{\beta ^2}{2\sqrt{m^2\mu ^2}}\left(C_1\mathrm{cos}\mu t+C_2\mathrm{sin}\mu t\right).$$
(121)
Including the zeroth order solution, the solution for $`q(t)`$ is given by
$$q^{(1)}(t)=C_1\left(\mathrm{cos}\mu t+\frac{\beta ^2}{4\mu \sqrt{m^2\mu ^2}}t\mathrm{sin}\mu t\right)+C_2\left(\mathrm{sin}\mu t\frac{\beta ^2}{4\mu \sqrt{m^2\mu ^2}}t\mathrm{cos}\mu t\right).$$
(122)
where we impose the initial condition so that $`q(0)=q^{(0)}(0)`$.
A problem is that the perturbation grows linearly in time. However, we need to be careful to interpret this growth of perturbations. In this toy model, we can easily find an exact solution. The corresponding exact solution becomes
$`q(t)`$ $`=`$ $`C_1\mathrm{cos}\left[\left(\mu {\displaystyle \frac{\beta ^2}{4\mu \sqrt{m^2\mu ^2}}}\right)t\right]+C_2\mathrm{sin}\left[\left(\mu {\displaystyle \frac{\beta ^2}{4\mu \sqrt{m^2\mu ^2}}}\right)t\right],`$ (123)
$`\varphi (y,t)`$ $`=`$ $`{\displaystyle \frac{\beta }{2\sqrt{m^2\mu ^2}}}q(t)e^{\sqrt{m^2\mu ^2}y},`$ (124)
for $`\beta 1`$. The effect of the coupling merely changes the frequency of the brane oscillator. The origin of the linear instability is that the naive expansion in terms of $`\beta `$ is not efficient. Indeed, we can use
$`\mathrm{cos}(A+B)`$ $`=`$ $`\mathrm{cos}A\mathrm{cos}B\mathrm{sin}A\mathrm{sin}B\mathrm{cos}AB\mathrm{sin}A,`$
$`\mathrm{sin}(A+B)`$ $`=`$ $`\mathrm{sin}A\mathrm{cos}B+\mathrm{sin}B\mathrm{cos}A\mathrm{sin}A+B\mathrm{cos}A,`$ (125)
for $`B1`$ and expand the exact solution into Eq.(122). However, this perturbation breaks down for large $`t`$. A crucial difference of the inflaton fluctuations case from the toy model is that the source term for the first order equation for $`q`$ contains a phase $`\phi `$ (compare Eqs. (109) and (110) to Eq. (121)). Then a perturbative solution cannot be written into the form like Eq. (123) using Eq. (125). This indicates that there could be a modification of the amplitude as well as the phase shift. The phase $`\phi `$ is originated from the fact that the brane oscillation cannot be matched by a single bound state (compare Eq. (100) and Eq. (120)). Thus this is an essential difference between the toy model and the inflaton fluctuations case.
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# Universal Approach to Overcoming Nonstationarity, Unsteadiness and Non-Markovity of Stochastic Processes in Complex Systems
## I Introduction
Nonstationarity, unsteadiness and non-Markovity are the most common essential peculiarities of stochastic processes in nature. The existence of the similar properties creates significant difficulties for the theoretical analysis of real complex systems Stanley . At present, methods, connected with localization of registered or calculated parameters for the quantitative account of the dramatic changes caused by the fast alternation of the behavior modes and intermittency, came into use. For example, the time behavior of the local (scale) Hurst exponents was defined in the recent work of Stanley et al. to study multifractal cascades in heartbeat dynamics Stanley and to analyze and forecast earthquakes and technogenic explosions in Ref. Yulmetyev1 . The application of the local characteristics allows one to avoid difficulties connected with nonergodicity of the investigated system and gives a possibility to extract additional valuable information on the hidden properties of real complex systems. From the physical point of view, this approach resembles the use of nonlinear equations of generalized hydrodynamics with the local time behavior of hydrodynamical and thermodynamical parameters and characteristics.
It is well known that one of the major problems of seismology is to predict the beginning of the main shock. Though science seems still to be far from the guaranteed decision of this problem there exist some interesting approaches based on the peculiar properties of precursory phenomena Sornette1 ; Igarashi ; Sornette2 ; Sornette3 ; Bufe ; Bak ; Sornette4 . Another important problem is recognition and differentiation of weak earthquakes and technogenic underground explosions signals. One of useful means of solving this problem is the defining of corresponding local characteristics Stanley ; Yulmetyev1 .
In the present work we suggest a new universal description of real complex systems by means of the microscopic, mesoscopic and macroscopic methods. We start with a macroscopic approach based on the kinetic theory of discrete stochastic processes and the hierarchy of the chain of finite-difference kinetic equations for the discrete time correlation function (TCF) and memory functions Yulmetyev1 ; Yulmetyev2 ; Yulmetyev3 .
The mesoscopic phenomena of the so-called โsoft matterโ physics, embracing a diverse range of system including liquid crystals, colloids, and biomembranes, generally involve some form of coupling of different characteristic time- and length-scales. Computational modelling of such multi-scale effects requires a new methodology applicable beyond the realm of traditional techniques such as ab initio and classical molecular dynamics (the methods of choice in the microscopic regime), as well as phase field modelling or the lattice-Boltzmann method (usually concerned with the macroscopic regime). As for complex systems, we propose to consider intermediate and slow processes within a unified framework of mesoscopic approach: by means of local time behavior of the local relaxation and kinetic parameters, local non-Markovity parameters and so on. For this purpose we introduce the notion of quasi-Brownian motion in a complex system by coarse-grained averaging of the initial time series on the basis of wavelet transformation.
As an example we consider here the local properties of relaxation or noise parameters for the analysis of seismic phenomena such as earthquakes and technogenic explosions. The layout of the paper is as follows. In Sec. II we describe in brief the stochastic dynamics of time correlation in complex systems by means of discrete non-Markov kinetic equations. Basic equations used for these calculations are presented here. The local noise parameters are defined in Sec. III. Section IV contains results obtained by the local noise parameter procedure for the case of seismic signals. The models of the time dependence of the local parameters are given in Sec. V. The basic conclusions are discussed in Sec. VI.
## II The basic definitions in kinetic description of discrete stochastic processes
### II.1 MACROSCOPIC DESCRIPTION IN THE ANALYSIS OF STOCHASTIC PROCESSES
A lot of different existent processes, such as economical, metheorogical, gravimetrical and other, are registered as discrete random series $`x_i`$ of some variable $`X`$. This random variable $`X`$ can be written as an array of its values
$$X=\{x(T),x(T+\tau ),x(T+2\tau ),\mathrm{},x(T+k\tau ),\mathrm{},x(T+(N1)\tau )\}.$$
(1)
Here a time step (or a time interval) $`\tau `$ is a constant, $`T`$ is the time when the registration of the signal begins, $`(N1)\tau `$ is the duration of the signal detection.
The average value $`x`$ and fluctuations $`\delta x_j`$ are defined by the following expressions, correspondingly
$$x=\frac{1}{N}\underset{j=0}{\overset{N1}{}}x(T+j\tau ),\delta x_j=\delta x(T+j\tau )=x(T+j\tau )x.$$
(2)
From the fluctuations of the considered random variable $`\delta x_j`$ we can form $`k`$-component state vector of the systemโs correlation state
$$\text{A}_k^0=\text{A}_k^0(0)=(\delta x(T),\delta x(T+\tau ),\mathrm{},\delta x(T+(k1)\tau ))=(\delta x_0,\delta x_1,\mathrm{},\delta x_{k1}).$$
(3)
The time dependence of the correlation state vector A can be represented as a discrete $`m`$-step shift
$$\text{A}_{m+k}^m=\text{A}_{m+k}^m(t)=(\delta x(T+m\tau ),\delta x(T+(m+1)\tau ),\mathrm{},\delta x(T+(m+k+1)\tau ))=(\delta x_0,\delta x_1,\mathrm{},\delta x_{m+k1}).$$
(4)
Then by analogy with the papers Yulmetyev1 ; Yulmetyev2 ; Yulmetyev3 we can write the following normalized TCF
$$M_0(t)=\frac{\text{A}_{N1m}^0\text{A}_{N1}^m}{\text{A}_{N1m}^0\text{A}_{N1m}^0}=\frac{\text{A}_{N1m}^0(0)\text{A}_{N1}^m(t)}{|\text{A}_{N1m}^0(0)|^2}=\frac{\text{A}_{N1m}^0(0)U(t=T+m\tau ,T)\text{A}_{N1m}^0(0)}{|\text{A}_{N1m}^0(0)|^2},$$
(5)
where angular brackets indicate the scalar product of the two state vectors. On the other hand, the time dependence of the vector $`\text{A}_{N1m}^m(T+t)`$, $`t=m\tau `$, can be represented formally with the help of the evolution operator $`U(t^{},t)`$ as follows:
$$\text{A}_{N1}^m(T+t)=U(T+m\tau ,T)\text{A}_{N1m}^0(T)=U(t,0)\text{A}_{N1m}^0(0).$$
(6)
The last one has the property: $`U(t,t)=1`$. Actually, one can write down formal discrete equation of motion with the use of the evolution operator $`U(t^{},t)`$ (see Appendix A for more details).
It was shown in Refs. Yulmetyev1 ; Yulmetyev2 ; Yulmetyev3 that the finite-difference kinetic equation of a non-Markov type for TCF $`M_0(t)`$ can be written by means of the technique of projection operators of Zwanzigโ-Moriโs type Zwanzig ; Mori as
$$\frac{\mathrm{}M_0(t)}{\mathrm{}t}=\lambda _1M_0(t)\tau \mathrm{\Lambda }_1\underset{j=0}{\overset{m1}{}}M_1(j\tau )M_0(tj\tau ).$$
(7)
Here the first order memory function $`M_1(j\tau )`$ appears, $`\lambda _1`$ is the eigenvalue of Liouvilleโs quasi-operator $`\widehat{}`$ and $`\mathrm{\Lambda }_1`$ is the relaxation noise parameter, which are characteristics of the investigated process. Possible methods of defining quasioperator $`\widehat{}`$ are presented in Appendix A \[see Eqs. (50),(52) and (54)\]. It should be recorded that Eq. (7) is the first kinetic finite-difference equation for initial TCF $`M_0(t)`$. With the use of the same procedure of projection operator we can obtain the chain of kinetic finite-difference equations of the following form:
$$\frac{\mathrm{}M_{i1}(t)}{\mathrm{}t}=\lambda _iM_{i1}(t)\tau \mathrm{\Lambda }_i\underset{j=0}{\overset{m1}{}}M_i(j\tau )M_{i1}(tj\tau ),i=1,2,3,\mathrm{}.$$
(8)
Here $`M_i(j\tau )`$ is the memory function of the $`i`$th order, whereas $`\lambda _i`$ and $`\mathrm{\Lambda }_i`$ are noise parameters:
$$\lambda _n=i\frac{\text{W}_{n1}\widehat{}\text{W}_{n1}}{|\text{W}_{n1}|^2};\mathrm{\Lambda }_n=i\frac{\text{W}_{n1}\widehat{}\text{W}_n}{|\text{W}_{n1}|^2}.$$
(9)
Here $`\text{W}_n`$ are the dynamical orthogonal variables, obtained by the Gram-Schmidt orthogonalization procedure
$$\text{W}_n,\text{W}_m=\delta _{n,m}|\text{W}_n|^2,$$
where $`\delta _{n,m}`$ is the Kroneckerโs symbol,
$$\text{W}_0=\text{A}_k^0(0),\text{W}_1=[i\widehat{}\lambda _1]\text{W}_0,$$
$$\text{W}_2=[i\widehat{}\lambda _2]\text{W}_1\mathrm{\Lambda }_1\text{W}_0,\mathrm{}.$$
(10)
From Eq. (10) it is obvious that in the cited Gram-Schmidt procedure from each new vector of state one should subtract the projection on all the previous vectors. Thereafter the orthogonalization (II.1) is complete.
A chain of integro-differential equations (8) arise as a result of the use of projection operator technique to define different correlation functions in physical problems Zwanzig ; Mori (for example, TCF of density fluctuation in Inelastic Neutron Scattering Yulmetyev\_Mokshin ; Scopigno and Light Scattering investigations, velocity autocorrelation function and others can be received in the specified way). However, the quest for the physically based way of closing the chain of equations and finding the spectra of the initial TCF are essential moments in these challenges. Here the situation is different. Namely, the initial TCF can be calculated directly from the experimental data. Then memory functions $`M_i(t)`$ and noise parameters $`\lambda _i`$, $`\mathrm{\Lambda }_i`$ are similarly calculated from the experiment. All these functions and parameters make it possible to carry out the detailed analysis of the random process. The noise parameters $`\lambda _i`$ and $`\mathrm{\Lambda }_i`$ are the relaxation characteristics of the experimental time series, which contain information of various modes passing and changing. The macroscopic approach presented above is based on the calculation of memory functions, power spectra, dynamic variables, relaxation parameters. It suggests the investigation of the system as a single whole. The global characteristics calculated on the basis of all the time series contain hidden information about various modes of the system behavior. As a rule, this information is difficult to extract and analyze. For this reason it is necessary to develop a mesoscopic description and introduce a local noise and relaxation parameters $`\lambda _i`$ and $`\mathrm{\Lambda }_i`$.
### II.2 MESOSCOPIC DESCRIPTION IN ANALYSIS OF STOCHASTIC PROCESSES
It is well known, the mesoscopic conception is one of the possible ways to deal with random processes in complex systems. It consists of extension of the domain of the dynamical equations mesoscopic variables and of introduction of a some local time interval. Such quantity as $`X`$ (it can be particle position, mass density and so on) is defined on the mesoscopic space mesoscop . In addition we introduce number $`M`$ as its extensive quantity. This number $`M`$ should satisfy the condition:
$$1MN,$$
(11)
where $`N`$ is the length of the initial experimental sampling. Let us take a working window of the fixed length $`M`$. By superposing this window on the initial sampling $`X`$, we choose all elements, incoming into it, as a separate sampling $`\overline{\xi }_0`$. Further, let us execute one time step $`\tau `$ shift of this working window to the right and obtain another local sampling of the length $`M`$. Executing this procedure ($`NM+1`$) times, one can obtain the same quantity of the local samplings of length $`M`$:
$`\overline{\xi }_0`$ $`=`$ $`\overline{\xi }_0\{x(T),x(T+\tau ),x(T+2\tau ),\mathrm{},x(T+(M1)\tau )\},`$
$`\overline{\xi }_1`$ $`=`$ $`\overline{\xi }_1\{x(T+\tau ),x(T+2\tau ),x(T+3\tau ),\mathrm{},x(T+M\tau )\},`$
$`\mathrm{},`$
$`\overline{\xi }_{NM}`$ $`=`$ $`\overline{\xi }_{NM}\{x(T+(NM1)\tau ),x(T+(NM)\tau ),x(T+(NM+1)\tau )\mathrm{},x(T+(N1)\tau )\}.`$ (12)
The obtained local samplings $`\overline{\xi }_i`$ form the array, which represents the time dynamics of the investigated process,
$$\overline{\xi }(t^{})=\{\overline{\xi }_0,\overline{\xi }_1,\overline{\xi }_2,\mathrm{},\overline{\xi }_i,\mathrm{},\overline{\xi }_{NM}\}.$$
(13)
Then, in accordance with the procedure described in the last subsection (II.1), we can define fluctuations by Eq. (2), calculate TCF with the help of Eq. (5), calculate memory functions and parameters $`\lambda _i`$ and $`\mathrm{\Lambda }_i`$ by Eqs. (9) for every sampling $`\overline{\xi }_i`$. However, parameters $`\lambda _i`$ and $`\mathrm{\Lambda }_i`$ will be characteristics of concrete $`j`$th sampling only. To characterize the local properties of the initial time data $`X`$, it is convenient to represent their local time dependence in the following way:
$`\lambda _i(t^{})=\{\lambda _i(T+(M1)\tau ),\lambda _i(T+M\tau ),\lambda _i(T+(M+1)\tau ),\mathrm{},\lambda _i(T+(NM)\tau )\},`$
$`\mathrm{\Lambda }_i(t^{})=\{\mathrm{\Lambda }_i(T+(M1)\tau ),\mathrm{\Lambda }_i(T+M\tau ),\mathrm{\Lambda }_i(T+(M+1)\tau ),\mathrm{},\mathrm{\Lambda }_i(T+(NM)\tau )\}.`$ (14)
Operating in the similar way, one can execute cross-over from the macroscopic description of the whole system to the mesoscopic one. The offered approach is very convenient for the description and analysis of non-stationary stochastic processes. It allows to depart from the global macro-characteristics, which carry only averaged minor information about the whole investigated process, and to turn to the stochastic description with the use of the local characteristics, and, as a result, to execute a more detailed analysis of various dynamic states of the system.
### II.3 CONCEPTION OF ONE-DIMENSIONAL DYNAMICS OF QUASI-BROWNIAN PARTICLE
Let us consider the motion of a large Brownian particle in a dense medium composed of light molecules and restrict it by a simple one-dimensional case. The coordinate $`x_i`$ and the velocity $`v_i`$ are random variables of a Brownian particle. The quantity $`\tau `$ represents average time between the two successive collisions of liquid molecules, $`T`$ is the initial moment of time. the variable $`M`$ characterizes a local size (mass) of a Brownian particle. It is obvious, that a Brownian particle must have a larger mass in comparison with liquid particles ($`M1`$), therefore it is more inert. Then it is convenient to define the coordinate of a Brownian particle at moment $`t^{}`$ as an average value of sampling $`\overline{\xi }_i`$. For example, we obtain from $`\overline{\xi }_0`$:
$$y_0=\frac{x(T)+x(T+\tau )+x(T+2\tau )+\mathrm{}+x(T+(M1)\tau )}{M}=\frac{1}{M}\underset{j=0}{\overset{M1}{}}x(T+j\tau ).$$
(15)
The quantity $`y_0`$ defines the coordinate of โthe center of massโ of a Brownian particle at the initial time moment $`t^{}=0`$. By analogy, it can define the position of a Brownian particle at the next time moment $`t^{}=\tau `$ and so on. As a result, we obtain a new time discrete series $`Y(t^{})`$ as:
$$Y(t^{})=\{y_0,y_1,y_2,\mathrm{},y_{NM}\}.$$
(16)
The velocity of a Brownian particle is
$$v_i=\frac{y_{i+1}y_i}{\tau }.$$
(17)
So, for example, for the initial velocity $`v_0`$ we obtain from Eqs. (12), (15) and (17):
$$v_0=\frac{y_1y_0}{\tau }=\frac{x(T+M\tau )x(T)}{M\tau }.$$
(18)
Obviously, if the larger one is $`M`$, the smaller one is the velocity of a Brownian particle. So, we obtain a discrete set of velocity values for a Brownian particle at equally small time interval $`\tau `$
$$V(t^{})=\{v_0,v_1,v_2,\mathrm{},v_{NM1}\}.$$
(19)
Eqs. (16) and (19) define time dependence of random variables $`Y(t^{})`$ and $`V(t^{})`$.
#### II.3.1 Generalized Langevin equation with discrete time
In accordance with Eq. (2) we define the correlation state vector, components of which are the fluctuations of the particle position,
$$\text{B}_k^0=\{\delta y_0,\delta y_1,\delta y_3,\mathrm{},\delta y_{k1}\}.$$
(20)
Then, by analogy to Eq. (4) time dependence of the vector B can be considered a discrete $`m`$-step time shift. For the vector B the following normalized TCF can be written with the help of Eq. (5)
$$b(t)=\frac{\text{B}_{N1m}^0\text{B}_{N1}^m}{|\text{B}_{N1m}^0|^2}.$$
(21)
The last one describes the time correlation of the two different correlation states of the system.
Now we introduce the linear projection operator in Euclidean space of the state vectors
$$Q\text{B}=\frac{\text{B}(0)\text{B}(0)\text{B}(t)}{|\text{B}(0)|^2}=\text{B}(0)b(t),Q=\frac{\text{B}(0)\text{B}(0)}{\text{B}(0)\text{B}(0)}.$$
(22)
This operator has the necessary property of idempotency $`Q^2=Q`$. The existence of projection operator $`Q`$ allows to introduce the mutually-supplementary projection operator $`R`$ as follows:
$$R=1Q,R^2=R,QR=RQ=0.$$
(23)
It is necessary to note that the projectors $`Q`$ and $`R`$ are both linear and can be recorded for the fulfillment of projection operations in the particular Euclidean space. The projection operators $`Q`$ and $`R`$ allow one to carry out the splitting of Euclidian space of vectors $`B`$, where $`B(0)`$ and $`B(t)B`$, into a straight sum of the two mutually supplementary subspaces in the following way:
$$B=B^{}+B^{\prime \prime },B^{}=QB,B^{\prime \prime }=RB.$$
(24)
Let us consider the finite-difference Liouvilleโs Eq. (48) for the vector of fluctuations of a Brownian particle position
$$\frac{\mathrm{}}{\mathrm{}t}\text{B}_{m+k}^m(t)=i\widehat{}(t,\tau )\text{B}_{m+k}^m(t).$$
(25)
Affecting the last equation by the operators $`Q`$ and $`R`$ successfully, we can split the dynamic equation (25) into two interconnected equations in the two mutually-supplementary Euclidean subspaces:
$`{\displaystyle \frac{\mathrm{}}{\mathrm{}t}}\text{B}^{}(t)=iQ\widehat{}(Q+R)\text{B}(t)=i\widehat{}_{11}\text{B}^{}(t)+i\widehat{}_{12}\text{B}^{\prime \prime }(t),`$ (26)
$`{\displaystyle \frac{\mathrm{}}{\mathrm{}t}}\text{B}^{\prime \prime }(t)=iR\widehat{}(Q+R)\text{B}(t)=i\widehat{}_{21}\text{B}^{}(t)+i\widehat{}_{22}\text{B}^{\prime \prime }(t).`$ (27)
In the above equations the matrix elements $`\widehat{}_{ij}`$ of the quasi-operator $`\widehat{}`$ have been introduced
$`\widehat{}_{11}=Q\widehat{}Q,\widehat{}_{12}=Q\widehat{}R,\widehat{}_{21}=R\widehat{}Q,\widehat{}_{22}=R\widehat{}R,`$
$`\widehat{}=\left(\begin{array}{cc}\widehat{}_{11}& \widehat{}_{12}\\ \widehat{}_{21}& \widehat{}_{22}\end{array}\right).`$ (28)
Operators $`\widehat{}_{ij}`$ act as follows: $`\widehat{}_{11}`$ \- from a subspaces $`B^{}`$ to $`B^{}`$; $`\widehat{}_{12}`$ \- from $`B^{\prime \prime }`$ to $`B^{}`$; $`\widehat{}_{21}`$ \- from $`B^{}`$ to $`B^{\prime \prime }`$; and $`\widehat{}_{22}`$ \- from $`B^{\prime \prime }`$ to $`B^{\prime \prime }`$.
To simplify the Liouville Eqs. (26) and (27), we exclude the irrelevant part $`B^{\prime \prime }(t)`$ and construct the closed equation for the relevant part $`B^{}(t)`$. For this purpose, it is necessary to obtain a step-by-step solution of Eq. (27)
$`{\displaystyle \frac{\mathrm{}\text{B}^{\prime \prime }(t)}{\mathrm{}t}}`$ $`=`$ $`{\displaystyle \frac{\text{B}^{\prime \prime }(t+\tau )\text{B}^{\prime \prime }(t)}{\tau }}=i\widehat{}_{21}\text{B}^{}(t)+i\widehat{}_{22}\text{B}^{\prime \prime }(t),`$ (29)
$`\text{B}^{\prime \prime }(t+\tau )`$ $`=`$ $`\text{B}^{\prime \prime }(t)+i\tau \widehat{}_{21}\text{B}^{}(t)+i\tau \widehat{}_{22}\text{B}^{\prime \prime }(t)`$ (30)
$`=`$ $`(1+i\tau \widehat{}_{22})\text{B}^{\prime \prime }(t)+i\tau \widehat{}_{21}\text{B}^{}(t)`$
$`=`$ $`U_{22}(t+\tau ,t)\text{B}^{\prime \prime }(t)+i\tau \widehat{}_{21}(t+\tau ,t)\text{B}^{}(t),`$
where $`U_{22}(t+\tau ,t)=1+i\tau \widehat{}_{22}(t+\tau ,t)`$ is the operator of a time step shift.
With the help of Eq. (30) we can derive the following expression for the next time step:
$`\text{B}^{\prime \prime }(t+2\tau )`$ $`=`$ $`U_{22}(t+2\tau ,t+\tau )\text{B}^{\prime \prime }(t+\tau )+i\tau \widehat{}_{21}(t+2\tau ,t+\tau )\text{B}^{}(t+\tau )`$ (31)
$`=`$ $`U_{22}(t+2\tau ,t+\tau )[U_{22}(t+\tau ,t)\text{B}^{\prime \prime }(t)+i\tau \widehat{}_{21}(t+\tau ,t)\text{B}^{}(t)]`$
$`+i\tau \widehat{}_{21}(t+2\tau ,t+\tau )\text{B}^{}(t+\tau )`$
$`=`$ $`U_{22}(t+2\tau ,t+\tau )U_{22}(t+\tau ,t)\text{B}^{\prime \prime }(t)+i\tau [U_{22}(t+2\tau ,t+\tau )`$
$`\times \widehat{}_{21}(t+\tau ,t)\text{B}^{}(t)+\widehat{}_{21}(t+2\tau ,t+\tau )\text{B}^{}(t+\tau )].`$
Generalizing this result in case of the $`m`$th discrete step we find the following final result
$`\text{B}^{\prime \prime }(t+m\tau )`$ $`=`$ $`\left\{\widehat{T}{\displaystyle \underset{j=0}{\overset{m1}{}}}U_{22}(t+(j+1)\tau ,t+j\tau )\right\}\text{B}^{\prime \prime }(t)`$ (32)
$`+i\tau {\displaystyle \underset{j=0}{\overset{m1}{}}}\left\{\widehat{T}{\displaystyle \underset{j^{}=j}{\overset{m2}{}}}U_{22}(t+(j^{}+2)\tau ,t+(j^{}+1)\tau )\right\}`$
$`\times \widehat{}(t+(j+1)\tau ,t+j\tau )\text{B}^{}(t+j\tau ).`$
Here $`\widehat{T}`$ denotes the Dyson operator of chronological ordering. Substituting the irrelevant part of Eq. (26) in the right side of Eq. (32), we obtain the closed finite-difference equation for the relevant part of the correlation state vector
$`{\displaystyle \frac{\mathrm{}}{\mathrm{}t}}\text{B}^{}(t+m\tau )`$ $`=`$ $`i\widehat{}_{11}(t+(m+1)\tau ,t+m\tau )\text{B}^{}(t+m\tau )+i\widehat{}_{12}(t+(m+1)\tau ,t+m\tau )`$ (33)
$`\times (\left\{\widehat{T}{\displaystyle \underset{j=0}{\overset{m1}{}}}U_{22}(t+(j+1)\tau ,t+j\tau )\right\}\text{B}^{\prime \prime }(t)\tau {\displaystyle \underset{j=0}{\overset{m1}{}}}\{\widehat{T}{\displaystyle \underset{j^{}=j}{\overset{m2}{}}}U_{22}(t+(j^{}+2)\tau ,t`$
$`+(j^{}+1)\tau )\}\widehat{}_{21}(t+(j+1)\tau ,t+j\tau )\text{B}^{}(t+j\tau )).`$
Substituting Eqs. (22) and (24) in Eq. (33), we derive a finite-difference kinetic equation of a non-Markov type for TCF $`b(t)`$
$`{\displaystyle \frac{\mathrm{}b(t)}{\mathrm{}t}}=\lambda _1b(t)\tau \mathrm{\Lambda }_1{\displaystyle \underset{j=0}{\overset{m1}{}}}M_1(tj\tau )b(j\tau ).`$ (34)
Here the TCF $`M_1(t)`$ is the first order memory function
$$M_1(tj\tau )=\frac{\text{W}_1(0)\text{W}_1(tj\tau )}{|\text{W}_1(0)|^2},$$
(35)
Here $`\mathrm{\Lambda }_1`$ is the frequency relaxation parameter of the first order with a square frequency dimension, and is defined by Eq. (9). Then if the dynamic variable $`\text{W}_0=\text{B}_k^0(0)`$ represents the fluctuations of a Brownian particle position, the dynamic variable $`\text{W}_1=i\widehat{}\text{W}_0\lambda _1\text{W}_0`$ contains fluctuations of pulses of a Brownian particle \[see Eqs. (10)\]. The function $`M_1(t)`$ is time correlation function of fluctuations of a Brownian particle velocity.
Defining the corresponding projection operators to new dynamic variable $`\text{W}_1`$ and repeating the above described procedure, we find the finite-difference kinetic equation for $`M_1(t)`$
$`{\displaystyle \frac{\mathrm{}M_1(t)}{\mathrm{}t}}=\lambda _2M_1(t)\tau \mathrm{\Lambda }_2{\displaystyle \underset{j=0}{\overset{m1}{}}}M_2(tj\tau )M_1(j\tau ).`$ (36)
Here $`\lambda _2`$ and $`\mathrm{\Lambda }_2`$ are the frequency relaxation parameters of the second order, $`M_2(t)`$ is the second order memory function or memory function of the velocity correlation function for a Brownian particle (memory friction) Tuckerman ; Berne ; Hanggi ; Hynes . In fact, Eq. (36) is a discrete finite-difference generalized Langevin equation (GLE). So, $`M_2(t)`$ can be associated with TCF of the stochastic Langevin forces, for which the similar equation can be received
$`{\displaystyle \frac{\mathrm{}M_2(t)}{\mathrm{}t}}=\lambda _3M_2(t)\tau \mathrm{\Lambda }_3{\displaystyle \underset{j=0}{\overset{m1}{}}}M_3(tj\tau )M_2(j\tau )`$ (37)
with the third order frequency relaxation parameters $`\lambda _3`$ and $`\mathrm{\Lambda }_3`$, and the memory function of the third order $`M_3(t)`$ respectively.
The three Eqs. (34), (36) and (37) are the exact consequence of microscopic discrete finite-difference equations of motion. Calculation of the memory function $`M_i(t)`$ and the relaxation frequency parameters $`\lambda _j`$, $`\mathrm{\Lambda }_k`$ are the central point here.
In case of a Brownian particle presented above, dynamic variables $`\text{W}_0`$ and $`\text{W}_1`$ are the position and pulses of random particles. In fact, they can be taken as characteristics of any other non-stationary process. The averaging operator of sampling with the length $`M`$, $`\widehat{A}=1/M_{j=0}^{M1}`$ \[see Eq. (15)\], applied to the intermediate local sampling, can be changed by any other operator, depending on the goal of the investigation. The operator $`\widehat{A}`$ allows one to get clear of sharp fluctuations in the initial sampling of data $`X(t)`$ and to replace it by an other $`Y(t^{})`$: $`X(t)Y(t^{})`$, which contains the results of coarse-graining averaging. However, it is not always convenient to average a local sampling. In these cases the operator $`\widehat{A}`$ can be replaced by an other operator, which allows one to obtain only one number from every local sampling. In general case, an another, a more universal method of transformation of the initial sampling into the sampling with some specified (required) characteristics can become discrete wavelet-transform Daub ; Wave\_1 ; Wave\_2 ; Wave\_3 ; Astafieva , which is defined by the following equation:
$$W(j,k)=\underset{j}{}\underset{k}{}X(k)2^{j2}\mathrm{\Psi }(2^jnk).$$
(38)
Here $`X(k)`$ is the sampling (1) and $`\mathrm{\Psi }(t)`$ is a time function with fast decay called mother wavelet. The following analysis can be applied to the new transformed data $`W(j,k)`$ according to the algorithm described above. Namely, memory functions $`M_i(t)`$ and frequency relaxation parameters $`\lambda _i`$ and $`\mathrm{\Lambda }_i`$ can be calculated for the transformed data.
#### II.3.2 Analogue of Green-Kubo relation for diffusion coefficient for time discrete system
According to the theory of random walkers the mean-square displacements of a Brownian particle can be defined as
$$\mathrm{}y^2_t=_{\mathrm{}}^+\mathrm{}\mathrm{}y^2\mathrm{\Phi }_1(y,t)๐y=2Dt,(t\mathrm{}).$$
(39)
Here $`\mathrm{\Phi }_1(y,t)`$ is the density of probability of being of particle at point $`y`$ at time moment $`t`$. The Eq. (39) is a well known Einstein relation for continuous displacements.
According to Eq. (39) the displacement during time $`t`$ is
$$\mathrm{}y_t=y(t)y(0).$$
(40)
In case of the discrete time Eq. (40) can be rewritten in the form
$$\mathrm{}y(t)=\mathrm{}y(T+j\tau )=y(T+(j+m)\tau )y(T+j\tau ),t=m\tau ,m1.$$
(41)
Then the mean-square displacement of a Brownian particle is
$$\mathrm{}y^2(T+(j+m)\tau )=\frac{1}{Nm}\underset{j=0}{\overset{Nm1}{}}[y(T+(j+m)\tau )y(T+j\tau )]^2,N>m1,$$
(42)
where $`Nm`$ is the quantity of โpossible waysโ. The diffusion coefficient takes the following form
$$D=\frac{1}{2m(Nm)\tau }\underset{j=0}{\overset{Nm1}{}}[y(T+(j+m)\tau )y(T+j\tau )]^2,\text{at}N>m1(\text{or}t\mathrm{}).$$
(43)
Letโs consider separately the sum in the last expression in terms of the velocity \[see Eq. (17)\]
$`{\displaystyle \underset{j=0}{\overset{Nm1}{}}}[`$ $`y`$ $`(T+(j+m)\tau )y(T+j\tau )]^2`$ (44)
$`=`$ $`\underset{(Nm)\text{of square brackets}[\mathrm{}]}{\underset{}{[y(T+m\tau )y(T)]^2+[y(T+(1+m)\tau )y(T+\tau )]^2+\mathrm{}+[y(T+(N1)\tau )y(T+(Nm1)\tau )]^2}}`$
$`=`$ $`[y(T+m\tau )y(T+(m1)\tau )+y(T+(m1)\tau )\mathrm{}y(T)]^2+\mathrm{}`$
$`[y(T+(N1)\tau )y(T+(N2)\tau )+y(T+(N2)\tau )\mathrm{}y(T+(Nm1)\tau )]^2`$
$`=`$ $`\tau ^2[v(T+(m1)\tau )+v(T+(m2)\tau )+\mathrm{}+v(T)]^2+\mathrm{}`$
$`+`$ $`\tau ^2[v(T+(N2)\tau )+v(T+(N3)\tau )+\mathrm{}+v(T+(Nm1)\tau )]^2`$
$`=`$ $`\tau ^2{\displaystyle \underset{j=0}{\overset{Nm1}{}}}\left[{\displaystyle \underset{j=0}{\overset{m1+j}{}}}v(T+k\tau )\right]^2.`$
Then the expression for the diffusion coefficient (43) can be rewritten as:
$$D=\frac{\tau }{2m(Nm)}\underset{j=0}{\overset{Nm1}{}}\left[\underset{k=j}{\overset{m1+j}{}}v(T+k\tau )\right]^2.$$
(45)
Eq. (45) is the discrete finite-difference analog of the famous Green-Kubo relation for the diffusion coefficient. Given relation has been obtained from Einstein equation (39) for a discrete system. The asymptotic limit $`t\mathrm{}`$ can be replaced here by the similar condition $`N\mathrm{}`$, $`m\mathrm{}`$, $`N>m`$.
## III Local noisy parameters
Seismic data represent discrete random series, which is recording of displacements of the Earthโs surface. Therefore, we can use the above-stated formalism to analyze seismic data. In particular, the local dependence of various characteristics Stanley ; Yulmetyev1 can be serve as additional source of information on properties of objects. The noise parameters $`\lambda _i`$ and $`\mathrm{\Lambda }_i`$ are very sensitive to the presence of a nonrandom component in the sampling. The change of the character of the correlated noise and, the appearance of the additional signal in the sampling can cause the alternation of these parameters. So, the time behavior of the local parameters $`\lambda _i`$, $`\mathrm{\Lambda }_i`$ is important and informative for analyzing of seismic data.
The procedure of localization consists in the following. Let us assume that we have an array of data $`\{x_1,x_2,x_3,\mathrm{},x_M,\mathrm{},x_N\}`$ and take the initial sampling of the fixed length $`M`$. Then by passing through all array of values with the โwork windowโ of the fixed length $`M`$ we can calculated the time series of the noisy parameters $`\{\lambda _i(T,T+M\tau ),\lambda _i(T+\tau ,T+(M+1)\tau ),\mathrm{},\lambda _i(T+(NM1)\tau ,T+(N1)\tau )\}`$ and $`\{\mathrm{\Lambda }_i(T,T+M\tau ),\mathrm{\Lambda }_i(T+\tau ,T+(M+1)\tau ),\mathrm{},\mathrm{\Lambda }_i(T+(NM1)\tau ,T+(N1)\tau )\}`$
Obviously, it is inadmissible to use both very large intervals ($`M=`$1000 and more points) and very short intervals ($`M=`$50 points and less) for definition of local parameters $`\lambda _i(t)`$, $`\mathrm{\Lambda }_i(t)`$. In the first case the physical sense of the localization procedure is lost. On the other hand, it is impossible to carry out any plausible correlation analysis with small intervals because of gross errors. Therefore there is necessity for finding the optimum length of the initial local interval or quantity $`M`$.
To determine the optimal minimal local sampling we have used the data corresponding to the calm state of the Earth before the technogenic explosion (an underground nuclear explosion). The calculation procedure consists in the following. We have taken the interval of $`40`$ points as the starting point and have calculated all the low-order noise parameters $`\lambda _i`$ $`(i=1,2,3)`$ and $`\mathrm{\Lambda }_j`$ $`(j=1,2)`$ by Eqs. (9) and (10). Then the interval was consistently increased by unit time segment $`\tau `$ and the relaxation parameters $`\lambda _i`$, $`\mathrm{\Lambda }_j`$ were calculated every time at the increase of the interval. As a result of this procedure executed for the calm state of the Earth we have established that all parameters take โsteadyโ numerical values at the interval approximately equal to $`150`$ points and more (see Fig. 1 for more details). Namely, from Fig. 1 one can see, that the parameters $`\lambda _1`$ and $`\mathrm{\Lambda }_1`$ take minimal values in their absolute quantity at this length of interval, the amplitude of value fluctuations of the parameter $`\mathrm{\Lambda }_2`$ gets lower and level off. It is again the evidence of reduction of the noise influence. Fluctuations of the parameters $`\lambda _2`$ and $`\lambda _3`$ also decrease, and the parameters themselves take steady values starting with the sampling of length $`150`$ points. Applying this procedure to other data of the calm state of the Earth, we find the same behavior of noise parameters $`\lambda _i`$, $`\mathrm{\Lambda }_i`$ and detect the minimal interval of $`150`$ points again. So, we choose the interval of such length as being optimal for accumulation of local statistics Herst .
## IV Definition of relaxation parameters $`\lambda _i(t)`$ and $`\mathrm{\Lambda }_j(t)`$ for earthquakes and technogenic explosions data
It is well known that modern seismic devices allow one to derive different quantitative and qualitative data about the seismic state of the Earth. In this work we analyse three various weak local earthquakes (EQโs) in Jordan (1998) \[EQ(1), EQ(2), EQ(3)\], one strong earthquake in Turkey (summer, 1999) \[EQ(T)\] and three local underground technogenic explosions (TE) \[TE(1), TE(2), TE(3)\] with the length of registration from $`10000`$ to $`25000`$ points. In case of strong EQ its seismogram contains $`65000`$ registered points. All these experimental data were courteously given by the Laboratory of Geophysics and Seismology (Amman, Jordan). All data correspond to transverse seismic displacements. The real temporal step of digitization $`\tau `$ between the registered points of seismic activity has the following values, viz, $`\tau =0.02s`$ for the EQ(T), and $`\tau =0.01s`$ for all others cases.
We defined such characteristics as the maximal amplitude of signal fluctuations before โeventโ $`a_1`$, the maximal amplitude of signal fluctuations during โeventโ $`a_2`$, EQ (TE) power $`a_2/a_1`$, time interval till EQ (TE) $`T_0`$, continuance of โeventโ $`T_l`$ and finally the total time of signal recording $`T_{total}`$ directly from seismic data. These quantities are presented in Table I. They give a clear notion about the duration and power of investigated phenomena. One can see from Table I, approximately $`45005700`$ points are accounted for the visible part a wavelet. This number of the recorded points allows one to execute analysis by means of the local parameters $`\lambda _i(t)`$ and $`\mathrm{\Lambda }_j(t)`$.
The procedure of calculation of time-dependence for $`\lambda _i(t)`$ and $`\mathrm{\Lambda }_j(t)`$ was based on the following operations. An interval with $`M150`$ points is taken, and noisy parameters $`\lambda _i`$, $`\mathrm{\Lambda }_j`$ are calculated for it with the help of Eqs. (9), (10). Then the operation of โstepwise shift to the rightโ at the interval of the fixed length $`M`$ is executed, and parameters are computed again. These actions are executed, while the initial sampling $`X(t)`$ will not be finished. As a result we obtain the following dependencies $`\{\lambda _i(T,T+M\tau ),\lambda _i(T+\tau ,T+(M+1)\tau ),\mathrm{},\lambda _i(T+(NM1)\tau ,T+(N1)\tau )\}`$ and $`\{\mathrm{\Lambda }_i(T,T+M\tau ),\mathrm{\Lambda }_i(T+\tau ,T+(M+1)\tau ),\mathrm{},\mathrm{\Lambda }_i(T+(NM1)\tau ,T+(N1)\tau )\}`$. If the character of the noise in the investigated data change, some signal will appear or disappear, and it will be directly reflected in the behavior of the relaxation characteristics.
The results of the above described procedure for the case of EQ and TE data are shown in Figs. $`2`$, $`3`$ and Table II. However, in order to check up the optimized length of the local interval, we calculated local parameters $`\lambda _i(t)`$ and $`\mathrm{\Lambda }_j(t)`$ $`(i=1,2,3)`$, $`(j=1,2)`$ at the local sampling with the length $`M=100,200,250,300,350`$ and $`400`$ points. It turned out that a large number of various noises in the behavior of $`\lambda _i(t)`$ and $`\mathrm{\Lambda }_j(t)`$ are superimposed on the carrying trajectory at $`M=100`$ points (parameters has a gross errors). The line-shapes of $`\lambda _i(t)`$ and $`\mathrm{\Lambda }_j(t)`$ practically cease to change at the sampling the length $`M=150`$ and more, $`M=200,250,\mathrm{}`$. Once again it is evidence that the local interval with the length $`M=150`$ points is optimal for the analysis of strong, weak EQโs and TEโs.
The detailed analysis of the received results allows one to reveal the following features.
Weak EQโs and local underground TEโs (for TE(3) and EQ(3), see Figs. 2, 3):
1. All the parameters $`\lambda _i(t)`$ take only negative values ($`\lambda _i(t)<0`$), whereas the noisy parameters $`\mathrm{\Lambda }_j(t)`$ can be both positive and negative.
2. Noisy parameter $`\lambda _1(t)`$. The absolute magnitude $`|\lambda _1|`$ increases sharply in its amplitude by a factor approximately equal to $`413.3`$ during EQ (various for different EQโs), and then it returns to its initial state. the restoration time $`T_{\lambda _1}`$ and the duration of โeventโ $`T_l`$ are approximately equal for weak EQ, i. e. $`T_{\lambda _1}T_l`$.
The absolute magnitude $`|\lambda _1|`$ also exhibits an abrupt rise $`2.23.5`$ times higher for TEโs. However, it returns quickly to its normal level. The restoration time $`\lambda _1(t)`$ for TE is less than the duration of the โeventโ approximately by a factor of $`2.53`$.
3. Noisy parameter $`\lambda _2(t)`$. The parameter $`\lambda _2(t)`$ responds to the beginning of the โeventโ by an abrupt rise of its value. It increases in amplitude both for TEโs and for EQโs. The character of the noise changes during the โeventโ. The parameter responds to the power of the โeventโ.
4. Noisy parameter $`\lambda _3(t)`$. This parameter always fluctuates near its numerical value $`1`$, and has an abrupt rise at the beginning the โeventโ in the form of separated spikes (see Figs. 3 and 4). The parameter keenly responds to the noise changes.
5. Noisy parameter $`\mathrm{\Lambda }_1(t)`$. It fluctuates near zero before and after the โeventโ changing its sign at this time. The parameter increases sharply at the โeventโ and always (!) retains positive. Then it decays smoothly. Restoration time $`T_{\mathrm{\Lambda }_1}`$ and the duration of the โeventโ $`T_l`$ are approximately the same both for the EQ and TE. This parameter is very sensitive to the changes of the noise character. For example, the data analysis of the EQ$`(2)`$ shows that the separate burst of the amplitude values of $`\mathrm{\Lambda }_1(t)`$ appears for $`4000`$ points up to the beginning of the EQ, although such indicator was not visually observed in the initial seismic data param . It may be the evidence of high prognostic property of this parameter for EQโs forecasting.
6. Noisy parameter $`\mathrm{\Lambda }_2(t)`$. Noise changes of the parameter $`\mathrm{\Lambda }_2(t)`$ are observed during the โeventโ both during EQโs and TEโs. From Figs. $`2`$ and $`3`$ we can see that $`\mathrm{\Lambda }_2(t)`$ has a distinctive negative depression during the โeventโ.
Strong EQโs (for EQ(T) see Fig. 4) Turkish :
All parameters react keenly to the appearance of the signal in case of a strong EQ. Let us consider the behavior of the local parameters in this case in detail.
7. Noisy parameter $`\lambda _1(t)`$. The parameter demonstrates the presence of noise and takes negative values before the โeventโ \[Region $`I`$ in Fig. 4 f)\]. As the โeventโ approaches, the amplitude of fluctuations decreases, and oscillations turn into negligible fluctuations near zero value for II, III and IV Regions (for more detail, see Fig. 4).
8. Noisy parameter $`\lambda _2(t)`$. A noise is also observed in the behavior of this parameter before the โeventโ, and this parameter takes only negative values. However, its values decrease sharply in absolute magnitude and begin to take values near zero with the appearance of an EQ signal (Fig. 4 b).
9. Noisy parameter $`\lambda _3(t)`$. This parameter takes only negative values at all times $`t`$. The negligible noise appears before the โeventโ. As the โeventโ approaches the noise increases in amplitude by the factor $`35`$. The amplitude of oscillations begins to decrease in Region IV (see Fig. 4 c).
10. Noisy parameters $`\mathrm{\Lambda }_1(t)`$ and $`\mathrm{\Lambda }_2(t)`$. Parameters fluctuate, taking positive values mainly before the โeventโ (Region I). With the beginning of โeventโ values of parameter increase greatly in absolute magnitude approaching zero. At the beginning of Region II both parameters change their sign from positive to negative. Then for Regions II, III and IV we see only right line with $`\mathrm{\Lambda }_1(t),\mathrm{\Lambda }_2(t)0`$ on the scales of Figs. 4 d) and e). However, the parameter $`\mathrm{\Lambda }_2(t)`$ has faintly visible fluctuations for Region IV characterizing the final EQ phase \[see Figs. 4 d) and e)\].
So, all parameters are very sensitive to the approach of a strong EQ. A sharp change in their behavior is appreciable before strong fluctuations of the Earthโs surface far off $`10000`$ points ($`3.5`$ minutes).
## V Simple exponential model for time local behavior of noisy relaxation parameters $`\lambda _1(t)`$ and $`\mathrm{\Lambda }_1(t)`$
As can be seen from Figs. $`24`$ all relaxation parameters are very sensitive to the beginning of EQ and TE. The behavior of parameters $`\lambda _1(t)`$ and $`\mathrm{\Lambda }_1(t)`$ in weak local EQโs and local TEโs is of great interest. These parameters are initial in our calculations \[see Eqs. (9), (10)\]. The analysis has shown that parameters oscillate near average values $`\lambda _0`$ and $`\mathrm{\Lambda }_0`$, correspondingly, before and after the oscillations visible on seismograms. The results of the behavior of these parameters for TE(3) and EQ(3) are presented in Figs. $`2`$ a)-d) and Fig. $`3`$ a)- d). However, a sudden rise by factors $`\mathrm{\Delta }\lambda _0`$ and $`\mathrm{\Delta }\mathrm{\Lambda }_0`$ is always observed in the behavior of these parameters at the enhancement or the appearance of the signal (it is seen at the beginning of the EQ or the TE in seismogram data). Furthermore, the continuous attenuation occurs. Over all this time these parameters have a well-defined pronounced trend. Such behavior of $`\lambda _1(t)`$ and $`\mathrm{\Lambda }_1(t)`$ give us a possibility of modelling the time dependence of these parameters by some simple mathematical functions.
The fitting procedure showed that the time behavior of these parameters can be well approximated by the following simple time dependencies:
$$\lambda _1(t)=\lambda _0+\mathrm{\Delta }\lambda \text{exp}\left\{\frac{tT_0}{T_\lambda }\right\}H(tT_0),$$
(46)
$$\mathrm{\Lambda }_1(t)=\mathrm{\Lambda }_0+\mathrm{\Delta }\mathrm{\Lambda }\text{exp}\left\{\frac{tT_0}{T_\mathrm{\Lambda }}\right\}H(tT_0),$$
(47)
where $`H(t)`$ is the Heaviside function, $`T_\lambda `$ and $`T_\mathrm{\Lambda }`$ are the relaxation times of $`\lambda _1(t)`$ and $`\mathrm{\Lambda }_1(t)`$, correspondingly. The time $`T_0`$ is the same in Eqs. (46) and (47) for parameters $`\lambda _1(t)`$ and $`\mathrm{\Lambda }_1(t)`$ (see Table 1). The numerical values of the variables, included in Eqs. (46), (47) were defined for EQโs and TEโs by comparison of localization results with these equations (see Fig. $`5`$). Numerical values of parameters are presented in Table 2.
So, it proved that the restoration of these parameters to their steady values occurs according to the exponential law. As can be seen from Fig. $`5`$ this description best suits for $`\lambda _1(t)`$ and $`\mathrm{\Lambda }_1(t)`$ of weak EQโs. The foregoing estimations strengthen fully our resume 2 and 5 of Section IV.
The results presented in the last three columns of Table II might be of interest for readers. As the quantities $`\mathrm{}\lambda _0/\lambda _0`$ and $`\mathrm{}\mathrm{\Lambda }_0/\mathrm{\Lambda }_0`$ show the rise of value of the corresponding parameter at TE and EQ. Finally, the ratio between the โeventโ duration $`T_l`$ and the relaxation time $`T_\lambda `$ discovers a remarkable distinction between TEโs and weak EQโs.
## VI Conclusion
In this work universal method for investigating non-stationary, unsteady and non-Markov random processes in discrete systems is suggested. This universality is achieved by combining the opportunities of microscopic, mesoscopic and macroscopic descriptions for complex systems. This method allows one to find and to analyze fast, slow and super-slow processes. To investigate super-slow processes we propose to use the model of a โquasi-Brownian particleโ. The wavelet-transformation of the initial time series can be used for this purpose. This method helps to analyze and differentiate similar signals of different origin. Theoretical investigations have been realized by means of two methods supplementing each other: the statistical theory of discrete non-Markov stochastic processes Yulmetyev1 and the local noisy parameters. The application of the last method gives a possibility to study non-stationary and unsteady processes with alternation and superimposition of different modes. The correlation between the time scales characteristic of different modes may be different. However, at accurate realization the localization procedure allows one to separate the noise and the signals Yulmetyev1 ; Stanley , and to carry out their quantitative and informative analysis.
Another important advantage of this method is the possibility to operate it in โreal timeโ regime, i.e. it can be put into practice immediately at getting the data, that is of great practical value.
The developed approach has been tested for strong and weak EQโs data and nuclear underground TEโs. As a result we have obtained of the following results.
The time-behavior of the local relaxation parameters can be described by simple model relaxation functions. The temporal relaxation of parameters $`\lambda _1(t)`$ and $`\mathrm{\Lambda }_1(t)`$ in weak EQโs and TEโs after the beginning of the โeventโ occurs according to the exponential law. However, the restoration and duration of the events are practically the same in case of weak EQโs. The restoration time of parameters $`\lambda _1(t)`$ and $`\mathrm{\Lambda }_1(t)`$ in the case of TEโs differs noticeably from the duration of the event. So, this approach can be useful in recognition these two different seismic phenomena. From the analysis of strong EQ data one can see that the behavior of all parameters changes greatly long before the EQ. For example, such change for the EQ(T) presented here occurs $`3.5`$ minutes before the main event. This change of $`\lambda _i(t)`$ and $`\mathrm{\Lambda }_i(t)`$ obtained for strong EQโs opens the possibility for a more accurate registration of the beginning of changes of the parameters before the visual wavelet and real EQโs.
We are sure that the suggested method can be very useful for the study of a wide class of random discrete processes in real complex systems of live and of lifeless things: in cardiology, physiology, neurophysiology, biophysics of membranes and seismology.
## VII Acknowledgments
We gratefully acknowledge Prof. Raoul Nigmatullin for stimulating discussion and Dr.L.O.Svirina for technical assistance. This work supported in part by Russian Fund for Basic Research (Grant No 02-02-16146), Ministry of Education of RF (Grant No A03-2.9-336, No 02-3.1-538) and RHSF (No 03-06-00218a).
## Appendix A Three forms of the quasi-operator $`\widehat{}`$
Equations of motion of a random variable $`x`$ with the use of the derivative of the three different forms Korn are represented here.
Using evolution operator, we can write the equation of motion for a discrete case as following
$$dx(t)/dt=i\widehat{}x(t).$$
(48)
However, there is a possibility of application of the time derivative $`d/dt\mathrm{}/\mathrm{}t`$ in three different forms:
1. โRightโ derivative (with decurrent difference in numerator)
$$\frac{\mathrm{}x(t)}{\mathrm{}t}=\frac{x(t+\tau )x(t)}{\tau }=\frac{1}{\tau }U(t+\tau ,t)x(t)$$
(49)
with Liouvilleโs quasioperator of the following form:
$$\widehat{}(t,\tau )=\frac{i}{\tau }[U(t+\tau ,t)1];$$
(50)
2. โLeftโ derivative (with ascending difference in numerator)
$$\frac{\mathrm{}x(t)}{\mathrm{}t}=\frac{x(t)x(t\tau )}{\tau }=\frac{x(t)U^1(t,t\tau )x(t)}{\tau }=\frac{1}{\tau }[1U^1(t,t\tau )]x(t)$$
(51)
with the Liouvillian of the next form:
$$\widehat{}(t,\tau )=\frac{i}{\tau }[1U^1(t,t\tau )];$$
(52)
3. โCentralโ derivative (with central difference in numerator)
$`{\displaystyle \frac{\mathrm{}x(t)}{\mathrm{}t}}`$ $`=`$ $`{\displaystyle \frac{x(t+\tau )x(t\tau )}{2\tau }}={\displaystyle \frac{x(t+\tau )x(t)}{2\tau }}{\displaystyle \frac{x(t)x(t\tau )}{2\tau }}`$ (53)
$`=`$ $`{\displaystyle \frac{1}{2\tau }}[U(t+\tau ,t)U^1(t,t\tau )]x(t).`$
Then the quasioperator $`\widehat{}`$ takes the following form:
$$\widehat{}(t,\tau )=\frac{i}{2\tau }[U(t+\tau ,t)U^1(t,t\tau )].$$
(54)
In the calculations and the analysis presented in this work we have used the derivative of the first form \[see Eqs. (49), (50)\].
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# Quasi-periodic solutions of the equation ๐ฃ_{๐กโข๐ก}-๐ฃ_{๐ฅโข๐ฅ}+๐ฃยณ=๐โข(๐ฃ)
## 1. Introduction
This paper*Keywords:* Nonlinear Wave Equation, Infinite dimensional Hamiltonian Systems, Quasi-periodic solutions, Lyapunov-Schmidt reduction. *2000AMS Subject Classification:* 35L05, 35B15, 37K50. Supported by MURST within the PRIN 2004 โVariational methods and nonlinear differential equationsโ. deals with a class of one-dimensional completely resonant nonlinear wave equations of the type
$`\{\begin{array}{c}v_{tt}v_{xx}=v^3+f(v)\hfill \\ v(t,x)=v(t,x+2\pi ),(t,x)^2\hfill \end{array}`$ (3)
where $`f:`$ is analytic in a neighborhood of $`v=0`$ and $`f(v)=๐ช(v^4)`$ as $`v0`$.
In the recent paper , M. Procesi proved the existence of small-amplitude quasi-periodic solutions of (3) of the form
$$v(t,x)=u(\omega _1t+x,\omega _2tx),$$
(4)
where $`u`$ is an odd analytic function, $`2\pi `$-periodic in both its arguments, and the frequencies $`\omega _1,\omega _21`$ belong to a Cantor-like set of zero Lebesgue measure. It is assumed that $`f`$ is odd and $`f(v)=๐ช(v^5)`$, see Theorem 1 in .
These solutions $`v(t,x)`$ correspond โ at the first order โ to the superposition of two waves, traveling in opposite directions:
$$v(t,x)=\sqrt{\epsilon }[r(\omega _1t+x)+s(\omega _2tx)+h.o.t.]$$
where $`\omega _1,\omega _2=1+๐ช(\epsilon )`$.
Motivated by the previous result, we study in the present paper the existence of quasi-periodic solutions of (3) having a different form, namely
$$v(t,x)=u(\omega _1t+x,\omega _2t+x).$$
(5)
Moreover we do not assume $`f`$ to be odd.
First of all, we have to consider different frequencies than in . Precisely, the appropriate choice for the relationship between the amplitude $`\epsilon `$ and the frequencies $`\omega _1,\omega _2`$ turns out to be
$$\omega _1=1+\epsilon +b\epsilon ^2,\omega _2=1+b\epsilon ^2,$$
(6)
where $`b1/2,\epsilon 0`$. This choice leads to look for quasi-periodic solutions $`v(t,x)`$ of (3) of the form
$$v(t,x)=u(\epsilon t,(1+b\epsilon ^2)t+x),$$
(7)
where $`(b,\epsilon )^2`$, $`\frac{1+b\epsilon ^2}{\epsilon }`$. On the contrary, taking in (5) frequencies $`\omega _1=1+\epsilon ,\omega _2=1+a\epsilon `$ as in , no quasi-periodic solutions can be found, see Remark in section 2. We show that there is no loss of generality passing from (5) to (7), because all the possible quasi-periodic solutions of (3) of the form (5) are of the form (7), see Appendix B.
Searching small amplitude quasi-periodic solutions of the form (7) by means of the Lyapunov-Schmidt method, leads to the usual system of a range equation and a bifurcation equation.
The former is solved, in a similar way as , by means of the standard Contraction Mapping Theorem, for a set of zero measure of the parameters. These arguments are carried out in section 4.
In section 5 we study the bifurcation equation, which is infinite-dimensional because we deal with a completely resonant equation. Here new difficulties have to be overcome. Since $`f`$ is not supposed to be odd, we cannot search odd solutions as in , so we look for even solutions. In this way, the bifurcation equation contains a new scalar equation for the average of $`u`$, see \[$`C`$-equation\] in (33), and the other equations contain supplementary terms.
To solve the bifurcation equation we use an ODE analysis; we cannot directly use variational methods as in ,, because we have to ensure that both components $`r,s`$ in (33) are non-trivial, in order to prove that the solution $`v`$ is actually quasi-periodic.
First, we find an explicit solution of the bifurcation equation (Lemma 1) by means of Jacobi elliptic functions (following ,,).
Next we prove its non-degeneracy (Lemmas 2,3,4); these computations are the heart of the present work. Instead of using a computer assisted proof as in , we here employ purely analytic arguments, see also (however, our problem requires much more involved computations than in ). In this way we prove the existence of quasi-periodic solutions of (3) of the form (7), see Theorem 1 (end of section 5).
ยฟFrom the physical point of view, this new class of solutions turns out to be, at the first order, the superposition of a traveling wave (with velocity greater than 1) and a modulation of long period, depending only on time:
$$v(t,x)=\epsilon [r(\epsilon t)+s((1+b\epsilon ^2)t+x)+h.o.t.].$$
Finally, in section 6 we show that our arguments can be also used to extend Procesi result to non-odd nonlinearities, see Theorem 2.
We also mention that recently existence of quasi-periodic solutions with $`n`$ frequencies have been proved in . The solutions found in belong to a neighborhood of a solution $`u_0(t)`$ periodic in time, independent of $`x`$, so they are different from the ones found in the present paper.
Acknowledgments. We warmly thank Massimiliano Berti for his daily support, Michela Procesi and Simone Paleari for some useful discussions.
## 2. The functional setting
We consider nonlinear wave equation (3),
$`\{\begin{array}{c}v_{tt}v_{xx}=v^3+f(v)\hfill \\ v(t,x)=v(t,x+2\pi )\hfill \end{array}`$
where $`f`$ is analytic in a neighborhood of $`v=0`$ and $`f(v)=๐ช(v^4)`$ as $`v0`$. We look for solutions of the form (5),
$$v(t,x)=u(\omega _1t+x,\omega _2t+x),$$
for $`(\omega _1,\omega _2)^2,`$$`\omega _1,\omega _21`$ and $`u`$$`2\pi `$-periodic in both its arguments. Solutions $`v(t,x)`$ of the form (5) are *quasi-periodic* in time $`t`$ when $`u`$ actually depends on both its arguments and the ratio between the periods is irrational, $`\frac{\omega _1}{\omega _2}`$.
We set the problem in the space $`_\sigma `$ defined as follows. Denote $`๐=/2\pi `$ the unitary circle, $`\phi =(\phi _1,\phi _2)๐^2`$. If $`u`$ is doubly $`2\pi `$-periodic, $`u:๐^2`$, its Fourier series is
$$u(\phi )=\underset{(m,n)^2}{}\widehat{u}_{mn}e^{im\phi _1}e^{in\phi _2}.$$
(9)
Let $`\sigma >0,s0`$. We define $`_\sigma `$ as the space of the even $`2\pi `$-periodic functions $`u:๐^2`$ which satisfy
$$\underset{(m,n)^2}{}\left|\widehat{u}_{mn}\right|^2\left[1+(m^2+n^2)^s\right]e^{2\sqrt{m^2+n^2}\sigma }:=u_\sigma ^2<\mathrm{}.$$
The elements of $`_\sigma `$ are even periodic functions which admit an analytic extension to the complex strip $`\{z:\left|\mathrm{Im}(z)\right|<\sigma \}`$.
$`(_\sigma ,_\sigma )`$ is a Hilbert space; for $`s>1`$ it is also an algebra, that is, there exists a constant $`c>0`$ such that
$$uv_\sigma cu_\sigma v_\sigma u,v_\sigma ,$$
see Appendix A. Moreover the inclusion $`_{\sigma ,s+1}_{\sigma ,s}`$ is compact.
We fix $`s>1`$ once and for all.
We note that all the possible quasi-periodic solutions of (3) of the form (5) are of the form (7) if we choose frequencies as in (6), see Appendix B. So we can look for solutions of (3) of the form (7), $`v(t,x)=u(\epsilon t,(1+b\epsilon ^2)t+x)`$, without loss of generality. For functions of the form (7), problem (3) is written as
$`\{\begin{array}{c}\epsilon \left[\epsilon _{\phi _1}^2+2(1+b\epsilon ^2)_{\phi _1\phi _2}^2+b\epsilon (2+b\epsilon ^2)_{\phi _2}^2\right](u)=u^3+f(u)\hfill \\ u_\sigma .\hfill \end{array}`$
We define $`M_{b,\epsilon }=\epsilon _{\phi _1}^2+2(1+b\epsilon ^2)_{\phi _1\phi _2}^2+b\epsilon (2+b\epsilon ^2)_{\phi _2}^2`$, rescale $`u\epsilon u`$ and set $`f_\epsilon (u)=\epsilon ^3f(\epsilon u)`$, so (3) can be written as
$`\{\begin{array}{c}M_{b,\epsilon }[u]=\epsilon u^3+\epsilon f_\epsilon (u)\hfill \\ u_\sigma .\hfill \end{array}`$ (13)
The main result of the present paper is the existence of solutions $`u_{(b,\epsilon )}`$ of (13) for $`(b,\epsilon )`$ in a suitable uncountable set (Theorem 1).
*Remark.* If we simply choose frequencies $`\omega _1=1+\epsilon ,\omega _2=1+a\epsilon `$ as in , we obtain a bifurcation equation different than (33). Precisely, it appears 0 instead of $`b(2+b\epsilon ^2)s^{\prime \prime }`$ in the left-hand term of the $`Q_2`$-equation in (33); so we do not find solutions which are non-trivial in both its arguments, but only solutions depending on the variable $`\phi _1`$. This is a problem because the quasi-periodicity condition requires dependence on both variables.
So we have to choose frequencies depending on $`\epsilon `$ in a more general way; a good choice is (6), $`\omega _1=1+\epsilon +b\epsilon ^2`$, $`\omega _2=1+b\epsilon ^2`$.
## 3. Lyapunov-Schmidt reduction
The operator $`M_{b,\epsilon }`$ is diagonal in the Fourier basis $`e_{mn}=e^{im\phi _1}e^{in\phi _2}`$ with eigenvalues $`D_{b,\epsilon }(m,n)`$, that is, if $`u`$ is written in Fourier series as in (9),
$`M_{b,\epsilon }[u]={\displaystyle \underset{(m,n)^2}{}}D_{b,\epsilon }(m,n)\widehat{u}_{mn}e^{im\phi _1}e^{in\phi _2},`$ (14)
where the eigenvalues $`D_{b,\epsilon }(m,n)`$ are given by
$`D_{b,\epsilon }(m,n)=`$ $`\epsilon m^2+2(1+b\epsilon ^2)mn+b\epsilon (2+b\epsilon ^2)n^2`$
$`=`$ $`(2+b\epsilon ^2)\left({\displaystyle \frac{\epsilon }{2+b\epsilon ^2}}m+n\right)\left(m+b\epsilon n\right).`$ (15)
For $`\epsilon =0`$ the operator is $`M_{b,0}=2_{\phi _1\phi _2}^2`$; its kernel $`Z`$ is the subspace of functions of the form $`u(\phi _1,\phi _2)=r(\phi _1)+s(\phi _2)`$ for some $`r,s_\sigma `$ one-variable functions,
$$Z=\{u_\sigma :\widehat{u}_{mn}=0(m,n)^2,m,n0\}.$$
We can decompose $`_\sigma `$ in four subspaces setting
$`\begin{array}{cc}& C=\{u_\sigma :u(\phi )=\widehat{u}_{0,0}\},\hfill \\ & Q_1=\{u_\sigma :u(\phi )=_{m0}\widehat{u}_{m,0}e^{im\phi _1}=r(\phi _1)\},\hfill \\ & Q_2=\{u_\sigma :u(\phi )=_{n0}\widehat{u}_{0,n}e^{in\phi _2}=s(\phi _2)\},\hfill \\ & P=\{u_\sigma :u(\phi )=_{m,n0}\widehat{u}_{mn}e^{im\phi _1}e^{in\phi _2}=p(\phi _1,\phi _2)\}.\hfill \end{array}`$ (20)
Thus the kernel is the direct sum $`Z=CQ_1Q_2`$ and the whole space is $`_\sigma =ZP`$. Any element $`u`$ can be decomposed as
$`\begin{array}{cc}\hfill u(\phi )=& \widehat{u}_{0,0}+r(\phi _1)+s(\phi _2)+p(\phi _1,\phi _2)\hfill \\ \hfill =& z(\phi )+p(\phi ).\hfill \end{array}`$ (23)
We denote $``$ the integral average: given $`g_\sigma `$,
$`\begin{array}{c}g=g_\phi =\frac{1}{(2\pi )^2}_0^{2\pi }_0^{2\pi }g(\phi )๐\phi _1๐\phi _2,\\ g_{\phi _1}=\frac{1}{2\pi }_0^{2\pi }g(\phi )๐\phi _1,g_{\phi _2}=\frac{1}{2\pi }_0^{2\pi }g(\phi )๐\phi _2.\end{array}`$
Note that $`\frac{1}{2\pi }_0^{2\pi }e^{ikt}๐t=0`$ for all integers $`k0,`$ so
$`\begin{array}{ccc}r=r_{\phi _1}=0\hfill & & r_{\phi _2}=r\hfill \\ s=s_{\phi _2}=0\hfill & & s_{\phi _1}=s\hfill \\ p=p_{\phi _1}=p_{\phi _2}=0\hfill & & u=\widehat{u}_{0,0}\hfill \end{array}`$
for all $`rQ_1,sQ_2,pP,u_\sigma ,`$ and by means of these averages we can construct the projections on the subspaces,
$`\mathrm{\Pi }_C=,\mathrm{\Pi }_{Q_1}=_{\phi _2},\mathrm{\Pi }_{Q_2}=_{\phi _1}.`$
Let $`u=z+p`$ as in (23); we write $`u^3`$ as $`u^3=z^3+(u^3z^3)`$ and compute the cube $`z^3=(\widehat{u}_{0,0}+r+s)^3`$. The operator $`M_{b,\epsilon }`$ maps every subspace of (20) in itself and it holds $`M_{b,\epsilon }[r]=\epsilon r^{\prime \prime }`$, $`M_{b,\epsilon }[s]=b\epsilon (2+b\epsilon ^2)s^{\prime \prime }`$, $`M_{b,\epsilon }[\widehat{u}_{0,0}]=0`$. So we can project our problem (13) on the four subspaces:
$`\begin{array}{cc}\hfill 0=& \widehat{u}_{0,0}^3+3\widehat{u}_{0,0}\left(r^2+s^2\right)+r^3+s^3+\hfill \\ & +\mathrm{\Pi }_C\left[(u^3z^3)f_\epsilon (u)\right]\left[C\text{-equation}\right]\hfill \\ \hfill r^{\prime \prime }=& 3\widehat{u}_{0,0}^2r+3\widehat{u}_{0,0}\left(r^2r^2\right)+r^3r^3+3s^2r+\hfill \\ & +\mathrm{\Pi }_{Q_1}\left[(u^3z^3)f_\epsilon (u)\right]\left[Q_1\text{-equation}\right]\hfill \\ \hfill b(2+b\epsilon ^2)s^{\prime \prime }=& 3\widehat{u}_{0,0}^2s+3\widehat{u}_{0,0}\left(s^2s^2\right)+s^3s^3+3r^2s+\hfill \\ & +\mathrm{\Pi }_{Q_2}\left[(u^3z^3)f_\epsilon (u)\right]\left[Q_2\text{-equation}\right]\hfill \\ \hfill M_{b,\epsilon }[p]=& \epsilon \mathrm{\Pi }_P\left[u^3+f_\epsilon (u)\right].\left[P\text{-equation}\right]\hfill \end{array}`$ (33)
Now we study separately the projected equations.
## 4. The range equation
We write the $`P`$-equation thinking $`p`$ as variable and $`z`$ as a โparameterโ,
$`M_{b,\epsilon }[p]=\epsilon \mathrm{\Pi }_P\left[(z+p)^3+f_\epsilon (z+p)\right].`$
We would like to invert the operator $`M_{b,\epsilon }`$. In Appendix C we prove that, fixed any $`\gamma (0,\frac{1}{4}),`$ there exists a non-empty uncountable set $`_\gamma ^2`$ such that, for all $`(b,\epsilon )_\gamma ,`$ it holds
$`\left|D_{b,\epsilon }(m,n)\right|>\gamma m,n,m,n0.`$
Precisely, our Cantor set $`_\gamma `$ is
$$_\gamma =\{(b,\epsilon )^2:\frac{\epsilon }{2+b\epsilon ^2},b\epsilon ^2\stackrel{~}{B}_\gamma ,\left|\frac{\epsilon }{2+b\epsilon ^2}\right|,|b\epsilon ^2|<\frac{1}{4},\frac{1+b\epsilon ^2}{\epsilon }\},$$
where $`\stackrel{~}{B}_\gamma `$ is a set of โbadly approximable numbersโ defined as
$`\stackrel{~}{B}_\gamma =\{x:\left|m+nx\right|>{\displaystyle \frac{\gamma }{|n|}}m,n,m0,n0\},`$ (34)
see Appendix C. Therefore $`M_{b,\epsilon |P}`$ is invertible for $`(b,\epsilon )_\gamma `$ and by (14) it follows
$`(M_{b,\epsilon |P})^1[h]={\displaystyle \underset{m,n0}{}}{\displaystyle \frac{\widehat{h}_{mn}}{D_{b,\epsilon }(m,n)}}e^{im\phi _1}e^{in\phi _2}`$
for every $`h=_{m,n0}\widehat{h}_{mn}e^{im\phi _1}e^{in\phi _2}P`$. Thus we obtain a bound for the inverse operators, uniformely in $`(b,\epsilon )_\gamma `$:
$`(M_{b,\epsilon |P})^1{\displaystyle \frac{1}{\gamma }}.`$
Applying the inverse operator $`(M_{b,\epsilon |P})^1`$, the $`P`$-equation becomes
$`p+\epsilon (M_{b,\epsilon |P})^1\mathrm{\Pi }_P\left[(z+p)^3f_\epsilon (z+p)\right]=0.`$ (35)
We would like to apply the Implicit Function Theorem, but the inverse operator $`(M_{b,\epsilon |P})^1`$ is defined only for $`(b,\epsilon )_\gamma `$ and in the set $`_\gamma `$ there are infinitely many holes, see Appendix C. So we fix $`(b,\epsilon )_\gamma ,`$ introduce an auxiliary parameter $`\mu `$ and consider the auxiliary equation
$`p+\mu (M_{b,\epsilon |P})^1\mathrm{\Pi }_P\left[(z+p)^3f_\mu (z+p)\right]=0.`$ (36)
Following Lemma 2.2 in , we can prove, by the standard Contraction Mapping Theorem, that there exists a positive constant $`c_1`$ depending only on $`f`$ such that, if
$$(\mu ,z)\times Z,|\mu |z_\sigma ^2<c_1\gamma ,$$
(37)
equation (36) admits a solution $`p_{(b,\epsilon )}(\mu ,z)P`$. Moreover, there exists a positive constant $`c_2`$ such that the solution $`p_{(b,\epsilon )}(\mu ,z)`$ respects the bound
$$p_{(b,\epsilon )}(\mu ,z)_\sigma \frac{c_2}{\gamma }z_\sigma ^3|\mu |.$$
(38)
Than we can apply the Implicit Function Theorem to the operator
$`\times Z\times PP`$
$`(\mu ,z,p)p+\mu (M_{b,\epsilon |P})^1\mathrm{\Pi }_P\left[(z+p)^3f_\mu (z+p)\right]`$
at every point $`(0,z,0)`$, so, by local uniqueness, we obtain the regularity: $`p_{(b,\epsilon )}`$, as function of $`(\mu ,z)`$, is at least of class $`๐^1`$.
Notice that the domain of any function $`p_{(b,\epsilon )}`$ is defined by (37), so it does not depend on $`(b,\epsilon )_\gamma `$.
In order to solve (35), we will need to evaluate $`p_{(b,\epsilon )}`$ at $`\mu =\epsilon `$; we will do it as last step, after the study of the bifurcation equation.
We observe that in these computations we have used the Hilbert algebra property of the space $`_\sigma ,uv_\sigma cu_\sigma v_\sigma u,v_\sigma .`$
## 5. The bifurcation equation
We consider auxiliary $`Z`$-equations: we put $`f_\mu `$ instead of $`f_\epsilon `$ in (33),
$`\begin{array}{cc}\hfill 0=& \widehat{u}_{0,0}^3+3\widehat{u}_{0,0}\left(r^2+s^2\right)+r^3+s^3+\hfill \\ & +\mathrm{\Pi }_C\left[(u^3z^3)f_\mu (u)\right]\left[Cequation\right]\hfill \\ \hfill r^{\prime \prime }=& 3\widehat{u}_{0,0}^2r+3\widehat{u}_{0,0}\left(r^2r^2\right)+r^3r^3+3s^2r+\hfill \\ & +\mathrm{\Pi }_{Q_1}\left[(u^3z^3)f_\mu (u)\right]\left[Q_1equation\right]\hfill \\ \hfill b(2+b\epsilon ^2)s^{\prime \prime }=& 3\widehat{u}_{0,0}^2s+3\widehat{u}_{0,0}\left(s^2s^2\right)+s^3s^3+3r^2s+\hfill \\ & +\mathrm{\Pi }_{Q_2}\left[(u^3z^3)f_\mu (u)\right].\left[Q_2equation\right]\hfill \end{array}`$ (45)
We substitute the solution $`p_{(b,\epsilon )}(\mu ,z)`$ of the auxiliary $`P`$-equation (36) inside the auxiliary $`Z`$-equations (45), writing $`u=z+p=z+p_{(b,\epsilon )}(\mu ,z)`$, for $`(\mu ,z)`$ in the domain (37) of $`p_{(b,\epsilon )}`$.
We have $`p_{(b,\epsilon )}(\mu ,z)=0`$ for $`\mu =0,`$ so the term $`\left[(u^3z^3)f_\mu (u)\right]`$ vanishes for $`\mu =0`$ and the bifurcation equations at $`\mu =0`$ become
$`\begin{array}{ccc}& \hfill 0=& \widehat{u}_{0,0}^3+3\widehat{u}_{0,0}\left(r^2+s^2\right)+r^3+s^3\hfill \\ & \hfill r^{\prime \prime }=& 3\widehat{u}_{0,0}^2r+3\widehat{u}_{0,0}\left(r^2r^2\right)+r^3r^3+3s^2r\hfill \\ & \hfill b(2+b\epsilon ^2)s^{\prime \prime }=& 3\widehat{u}_{0,0}^2s+3\widehat{u}_{0,0}\left(s^2s^2\right)+s^3s^3+3r^2s.\hfill \end{array}`$ (49)
We look for non-trivial $`z=\widehat{u}_{0,0}+r(\phi _1)+s(\phi _2)`$ solution of (49). We rescale setting
$`\begin{array}{cc}r=x\hfill & \widehat{u}_{0,0}=c\hfill \\ s=\sqrt{b(2+b\epsilon ^2)}y\hfill & \lambda =\lambda _{b,\epsilon }=b(2+b\epsilon ^2),\hfill \end{array}`$ (52)
so the equations become
$`\begin{array}{cc}& c^3+3c\left(x^2+\lambda y^2\right)+x^3+\lambda ^{3/2}y^3=0\hfill \\ & x^{\prime \prime }+3c^2x+3c\left(x^2x^2\right)+x^3x^3+3\lambda y^2x=0\hfill \\ & y^{\prime \prime }+3c^2\frac{1}{\lambda }y+3c\frac{1}{\sqrt{\lambda }}\left(y^2y^2\right)+y^3y^3+3\frac{1}{\lambda }x^2y=0.\hfill \end{array}`$ (56)
In the following we show that, for $`|\lambda 1|`$ sufficiently small, the system (56) admits a non-trivial non-degenerate solution. We consider $`\lambda `$ as a free real parameter, recall that $`Z=C\times Q_1\times Q_2`$ and define $`G:\times ZZ`$ setting $`G(\lambda ,c,x,y)`$ as the set of three left-hand terms of (56).
Lemma 1. *There exist $`\overline{\sigma }>0`$ and a non-trivial one-variable even analytic function $`\beta _0`$ belonging to $`_\sigma `$ for every $`\sigma (0,\overline{\sigma })`$, such that $`G(1,0,\beta _0,\beta _0)=0`$, that is $`(0,\beta _0,\beta _0)`$ solves (56) for $`\lambda =1`$.*
*Proof.* We prove the existence of a non-trivial even analytic function $`\beta _0`$ which satisfies
$`\beta _0^{\prime \prime }+\beta _0^3+3\beta _0^2\beta _0=0,\beta _0=\beta _0^3=0.`$ (57)
For any $`m(0,1)`$ we consider the Jacobi amplitude $`\mathrm{am}(,m):`$ as the inverse of the elliptic integral of the first kind
$$I(,m):,I(\phi ,m)=_0^\phi \frac{d\vartheta }{\sqrt{1m\mathrm{sin}^2\vartheta }}.$$
We define the Jacobi elliptic cosine setting
$$\mathrm{cn}(\xi )=\mathrm{cn}(\xi ,m)=\mathrm{cos}(\mathrm{am}(\xi ,m)),$$
see ch.16, . It is a periodic function of period $`4K`$, where $`K=K(m)`$ is the complete elliptic integral of the first kind
$$K(m)=_0^{\pi /2}\frac{d\vartheta }{\sqrt{1m\mathrm{sin}^2\vartheta }}.$$
Jacobi cosine is even, and it is also odd-symmetric with respect to $`K`$ on $`[0,2K],`$ that is $`\mathrm{cn}(\xi +K)=\mathrm{cn}(\xi K)`$, just like the usual cosine. Then the averages on the period $`4K`$ are
$`\mathrm{cn}=\mathrm{cn}^3=0.`$
Therefore it admits an analytic extension with a pole at $`iK^{}`$, where $`K^{}=K(1m)`$, and it satisfies $`(\mathrm{cn}^{})^2=m\mathrm{cn}^4+(2m1)\mathrm{cn}^2+(1m),`$ then $`\mathrm{cn}`$ is a solution of the ODE
$`\mathrm{cn}^{\prime \prime }+2m\mathrm{cn}^3+(12m)\mathrm{cn}=0.`$
We set $`\beta _0(\xi )=V\mathrm{cn}(\mathrm{\Omega }\xi ,m)`$ for some real parameters $`V,\mathrm{\Omega }>0,m(0,1)`$. $`\beta _0`$ has a pole at $`i\frac{K^{}}{\mathrm{\Omega }}`$, so it belongs to $`_\sigma `$ for every $`0<\sigma <\frac{K^{}}{\mathrm{\Omega }}`$.
$`\beta _0`$ satisfies
$`\beta _0^{\prime \prime }+\left(2m{\displaystyle \frac{\mathrm{\Omega }^2}{V^2}}\right)\beta _0^3+\mathrm{\Omega }^2(12m)\beta _0=0.`$
If there holds the equality $`\mathrm{\hspace{0.17em}2}m\mathrm{\Omega }^2=V^2,`$ the equation becomes
$`\beta _0^{\prime \prime }+\beta _0^3+\mathrm{\Omega }^2(12m)\beta _0=0.`$
$`\beta _0`$ is $`\frac{4K(m)}{\mathrm{\Omega }}`$-periodic; it is $`2\pi `$-periodic if $`\mathrm{\Omega }=\frac{2K(m)}{\pi }`$. Hence we require
$`2m\mathrm{\Omega }^2=V^2,\mathrm{\Omega }={\displaystyle \frac{2K(m)}{\pi }}.`$ (58)
The other Jacobi elliptic functions we will use are
$$\mathrm{sn}(\xi )=\mathrm{sin}(\mathrm{am}(\xi ,m)),\mathrm{dn}(\xi )=\sqrt{1m\mathrm{sn}^2(\xi )},$$
see ,. From the equality $`m\mathrm{cn}^2(\xi )=\mathrm{dn}^2(\xi )(1m),`$ with change of variable $`x=\mathrm{am}(\xi )`$ we obtain
$$_0^{K(m)}m\mathrm{cn}^2(\xi )๐\xi =E(m)(1m)K(m),$$
where $`E(m)`$ is the complete elliptic integral of the second kind,
$$E(m)=_0^{\pi /2}\sqrt{1m\mathrm{sin}^2\vartheta }๐\vartheta .$$
Thus the average on $`[0,2\pi ]`$ of $`\beta _0^2`$ is
$`\beta _0^2={\displaystyle \frac{V^2}{mK(m)}}\left[E(m)(1m)K(m)\right].`$
We want the equality $`\mathrm{\hspace{0.17em}3}\beta _0^2=\mathrm{\Omega }^2(12m)`$ and this is true if
$`E(m)+{\displaystyle \frac{8m7}{6}}K(m)=0.`$ (59)
The left-hand term $`\psi (m):=E(m)+\frac{8m7}{6}K(m)`$ is continuous in $`m`$; its value at $`m=0`$ is $`(\pi /12)<0`$, while at $`m=1/2`$, by definition of $`E`$ and $`K`$,
$$\psi \left(\frac{1}{2}\right)=\frac{1}{2}_0^{\pi /2}\frac{\mathrm{cos}^2\vartheta }{\left(1\frac{1}{2}\mathrm{sin}^2\vartheta \right)^{1/2}}๐\vartheta >0.$$
Moreover, its derivative is strictly positive for every $`m[0,\frac{1}{2}]`$,
$$\begin{array}{cc}\hfill \psi ^{}(m)=& _0^{\pi /2}\frac{8\frac{5}{2}\mathrm{sin}^2\vartheta +3m\mathrm{sin}^4\vartheta 8m\mathrm{sin}^2\vartheta }{6(1m\mathrm{sin}^2\vartheta )^{3/2}}๐\vartheta \hfill \\ \hfill & _0^{\pi /2}\frac{3+\mathrm{cos}^2\vartheta }{6}๐\vartheta >0,\hfill \end{array}$$
hence there exists a unique $`\overline{m}(0,\frac{1}{2})`$ which solves (59). Thanks to the tables in , p. 608-609, we have $`0.20<\overline{m}<0.21`$.
By (58) the value $`\overline{m}`$ determines the parameters $`\overline{\mathrm{\Omega }}`$ and $`\overline{V}`$, so the function $`\beta _0(\xi )=\overline{V}\mathrm{cn}(\overline{\mathrm{\Omega }}\xi ,\overline{m})`$ satisfies (57) and $`(0,\beta _0,\beta _0)`$ is a solution of (56) for $`\lambda =1`$. Therefore $`\beta _0_\sigma `$ for every $`\sigma (0,\overline{\sigma })`$, where $`\overline{\sigma }=(\frac{K^{}}{\mathrm{\Omega }})_{|m=\overline{m}}.`$ $`\mathrm{}`$
The next step will be to prove the non-degeneracy of the solution $`(1,0,\beta _0,\beta _0)`$, that is to show that the partial derivative $`_ZG(1,0,\beta _0,\beta _0)`$ is an invertible operator. This is the heart of the present paper. We need some preliminary results.
Lemma 2. *Given $`h`$ even $`2\pi `$-periodic, there exists a unique even $`2\pi `$-periodic $`w`$ such that*
$$w^{\prime \prime }+\left(3\beta _0^2+3\beta _0^2\right)w=h.$$
*This defines the Green operator $`L:_\sigma _\sigma `$, $`L[h]=w`$.*
*Proof.* We fix a $`2\pi `$-periodic even function $`h`$. We look for even $`2\pi `$-periodic solutions of the non-homogeneous equation
$`x^{\prime \prime }+\left(3\beta _0^2+3\beta _0^2\right)x=h.`$ (60)
First of all, we construct two solutions of the homogeneous equation
$`x^{\prime \prime }+\left(3\beta _0^2+3\beta _0^2\right)x=0.`$ (61)
We recall that $`\beta _0`$ satisfies $`\beta _0^{\prime \prime }+\beta _0^3+3\beta _0^2\beta _0=0`$, then deriving with respect to its argument $`\xi `$ we obtain $`\beta _0^{\prime \prime \prime }+3\beta _0^2\beta _0^{}+3\beta _0^2\beta _0^{}=0,`$ so $`\beta _0^{}`$ satisfies (61). We set
$`\overline{u}(\xi )={\displaystyle \frac{1}{\overline{V}\overline{\mathrm{\Omega }}^2}}\beta _0^{}(\xi )={\displaystyle \frac{1}{\overline{\mathrm{\Omega }}}}\mathrm{cn}^{}(\overline{\mathrm{\Omega }}\xi ,\overline{m}),`$ (62)
thus $`\overline{u}`$ is the solution of the homogeneous equation such that $`\overline{u}(0)=0`$, $`\overline{u}^{}(0)=1`$. It is odd and $`2\pi `$-periodic.
Now we construct the other solution. We indicate $`c_0`$ the constant $`c_0=\beta _0`$. We recall that, for any $`V,\mathrm{\Omega },m`$ the function $`y(\xi )=V\mathrm{cn}(\mathrm{\Omega }\xi ,m)`$ satisfies
$`y^{\prime \prime }+\left(2m{\displaystyle \frac{\mathrm{\Omega }^2}{V^2}}\right)y^3+\mathrm{\Omega }^2(12m)y=0.`$
We consider $`m`$ and $`V`$ as functions of the parameter $`\mathrm{\Omega }`$, setting
$$m=m(\mathrm{\Omega })=\frac{1}{2}\frac{3c_0}{2\mathrm{\Omega }^2},V=V(\mathrm{\Omega })=\sqrt{\mathrm{\Omega }^23c_0}.$$
(63)
We indicate $`y_\mathrm{\Omega }(\xi )=V(\mathrm{\Omega })\mathrm{cn}(\mathrm{\Omega }\xi ,m(\mathrm{\Omega })),`$ so $`(y_\mathrm{\Omega })_\mathrm{\Omega }`$ is a one-parameter family of solutions of
$$y_\mathrm{\Omega }^{\prime \prime }+y_\mathrm{\Omega }^3+3c_0y_\mathrm{\Omega }=0.$$
We can derive this equation with respect to $`\mathrm{\Omega }`$, obtaining
$$(_\mathrm{\Omega }y_\mathrm{\Omega })^{\prime \prime }+3y_\mathrm{\Omega }^2(_\mathrm{\Omega }y_\mathrm{\Omega })+3c_0(_\mathrm{\Omega }y_\mathrm{\Omega })=0.$$
Now we evaluate $`(_\mathrm{\Omega }y_\mathrm{\Omega })`$ at $`\mathrm{\Omega }=\overline{\mathrm{\Omega }}`$, where $`\overline{\mathrm{\Omega }}`$ correspond to the value $`\overline{m}`$ found in Lemma 1. For $`\mathrm{\Omega }=\overline{\mathrm{\Omega }}`$ it holds $`y_{\overline{\mathrm{\Omega }}}=\beta _0`$, so $`(_\mathrm{\Omega }y_\mathrm{\Omega })_{|\mathrm{\Omega }=\overline{\mathrm{\Omega }}}`$ satisfy (61). In order to normalize this solution, we compute
$$(_\mathrm{\Omega }y_\mathrm{\Omega })(\xi )=(_\mathrm{\Omega }V)\mathrm{cn}(\mathrm{\Omega }\xi ,m)+V\xi \mathrm{cn}^{}(\mathrm{\Omega }\xi ,m)+V_m\mathrm{cn}(\mathrm{\Omega }\xi ,m)(_\mathrm{\Omega }m).$$
Since $`\mathrm{cn}(0,m)=1`$$`m`$, it holds $`_m\mathrm{cn}(0,m)=0`$; therefore $`\mathrm{cn}^{}(0,m)=0`$$`m`$. From (63) we have $`_\mathrm{\Omega }V=\frac{\mathrm{\Omega }}{V}`$, so we can normalize setting
$$\overline{v}(\xi )=\frac{\overline{V}}{\overline{\mathrm{\Omega }}}\left(_\mathrm{\Omega }y_\mathrm{\Omega }\right)_{|\mathrm{\Omega }=\overline{\mathrm{\Omega }}}(\xi ).$$
$`\overline{v}`$ is the solution of the homogeneous equation (61) such that $`\overline{v}(0)=1`$, $`\overline{v}^{}(0)=0`$. We can write an explicit formula for $`\overline{v}`$. From the definitions it follows for any $`m`$
$$_m\mathrm{am}(\xi ,m)=\mathrm{dn}(\xi ,m)\frac{1}{2}_0^\xi \frac{\mathrm{sn}^2(t,m)}{\mathrm{dn}^2(t,m)}๐t.$$
Therefore $`\mathrm{cn}^{}(\xi )=\mathrm{sn}(\xi )\mathrm{dn}(\xi )`$; then we obtain for $`(V,\mathrm{\Omega },m)=(\overline{V},\overline{\mathrm{\Omega }},\overline{m})`$
$`\overline{v}(\xi )=\mathrm{cn}(\overline{\mathrm{\Omega }}\xi )+{\displaystyle \frac{\overline{V}^2}{\overline{\mathrm{\Omega }}}}\mathrm{cn}^{}(\overline{\mathrm{\Omega }}\xi )\left[\xi +{\displaystyle \frac{2\overline{m}1}{2}}{\displaystyle _0^\xi }{\displaystyle \frac{\mathrm{sn}^2(\overline{\mathrm{\Omega }}t)}{\mathrm{dn}^2(\overline{\mathrm{\Omega }}t)}}๐t\right].`$ (64)
By formula (64) we can see that $`\overline{v}`$ is even; it is not periodic and there holds
$`\overline{v}(\xi +2\pi )\overline{v}(\xi )={\displaystyle \frac{\overline{V}^2k}{\overline{\mathrm{\Omega }}}}\mathrm{cn}^{}(\overline{\mathrm{\Omega }}\xi )=\overline{V}^2k\overline{u}(\xi ),`$ (65)
where
$$k:=\mathrm{\hspace{0.17em}2}\pi +\frac{2\overline{m}1}{2}_0^{2\pi }\frac{\mathrm{sn}^2(\overline{\mathrm{\Omega }}t)}{\mathrm{dn}^2(\overline{\mathrm{\Omega }}t)}๐t.$$
(66)
ยฟFrom the equalities (L.1) and (L.2) of Lemma 3 we obtain
$$k=2\pi \frac{1+16\overline{m}16\overline{m}^2}{12\overline{m}(1\overline{m})},$$
(67)
so $`k>0`$ because $`\overline{m}(0.20,0.21)`$.
We have constructed two solutions $`\overline{u},\overline{v}`$ of the homogeneous equation; their wronskian $`\overline{u}^{}\overline{v}\overline{u}\overline{v}^{}`$ is equal to $`1`$, so we can write a particular solution $`\overline{w}`$ of the non-homogeneous equation (60) as
$`\overline{w}(\xi )=\left({\displaystyle _0^\xi }h\overline{v}\right)\overline{u}(\xi )\left({\displaystyle _0^\xi }h\overline{u}\right)\overline{v}(\xi ).`$
Every solution of (60) is of the form $`w=A\overline{u}+B\overline{v}+\overline{w}`$ for some $`(A,B)^2`$. Since $`h`$ is even, $`\overline{w}`$ is also even, so $`w`$ is even if and only if $`A=0`$.
An even function $`w=B\overline{v}+\overline{w}`$ is $`2\pi `$-periodic if and only if $`w(\xi +2\pi )w(\xi )=0`$, that is, by (65),
$$\left(_\xi ^{\xi +2\pi }h\overline{v}\right)\overline{u}(\xi )+\left[\left(_0^\xi h\overline{u}\right)B\right]\overline{V}^2k\overline{u}(\xi )=0\xi .$$
We remove $`\overline{u}(\xi )`$, derive the expression with respect to $`\xi `$ and from (65) it results zero at any $`\xi `$. Then the expression is a constant; we compute it at $`\xi =0`$ and obtain, since $`h\overline{u}`$ is odd and $`2\pi `$-periodic, that $`w`$ is $`2\pi `$-periodic if and only if $`B=\frac{1}{\overline{V}^2k}_0^{2\pi }h\overline{v}`$.
Thus, given $`h`$ even $`2\pi `$-periodic, there exists a unique even $`2\pi `$-periodic $`w`$ such that $`w^{\prime \prime }+\left(3\beta _0^2+3\beta _0^2\right)w=h`$ and this defines the operator $`L`$,
$`L[h]=\left({\displaystyle _0^\xi }h\overline{v}\right)\overline{u}(\xi )+\left[\left({\displaystyle \frac{1}{\overline{V}^2k}}{\displaystyle _0^{2\pi }}h\overline{v}\right){\displaystyle _0^\xi }h\overline{u}\right]\overline{v}(\xi ).`$ (68)
$`L`$ is linear and continuous with respect to $`_\sigma `$; it is the Green operator of the equation $`x^{\prime \prime }+\left(3\beta _0^2+3\beta _0^2\right)x=h`$, so, by classical arguments, it is a bounded operator of $`_{\sigma ,s}`$ into $`_{\sigma ,s+2}`$; the inclusion $`_{\sigma ,s+2}_{\sigma ,s}`$ is compact, then $`L:_\sigma _\sigma `$ is compact. $`\mathrm{}`$
Lemma 3. *There holds the following equalities and inequalities:*
(L.1) $`\mathrm{cn}^2=\frac{12\overline{m}}{6\overline{m}}`$ for $`m=\overline{m}`$. (Recall: $`\mathrm{cn}=\mathrm{cn}(,m))`$
(L.2) $`\frac{\mathrm{sn}^2}{\mathrm{dn}^2}=\frac{1}{1m}\mathrm{cn}^2`$ for any $`m`$.
(L.3) $`m\mathrm{cn}^2\frac{\mathrm{sn}^2}{\mathrm{dn}^2}=12\mathrm{cn}^2`$ for any $`m`$.
(L.4) *Exchange rule.* $`gL[h]=hL[g]g,h`$ even $`2\pi `$-periodic.
(L.5) $`13\beta _0^2L[1]=3\beta _0^2L[1]`$.
(L.6) $`\beta _0^2L[\beta _0]=\beta _0^2L[\beta _0]`$.
(L.7) $`3\beta _0^2L[\beta _0^2]=\beta _0^2\left(13L[\beta _0^2]\right)`$.
(L.8) $`\beta _0^2L[\beta _0]=\beta _0L[\beta _0^2]=L[\beta _0]=0`$.
(L.9) $`A_0:=13\beta _0^2L[1]0`$.
(L.10) $`B_0:=16\beta _0L[\beta _0]0`$ .
(L.11) $`C_0:=1+6\beta _0L[\beta _0]0`$ .
(L.12) $`A_01,L[\beta _0^2]0.`$
*Proof.* (L.1) By construction of $`\beta _0`$ we have $`\overline{\mathrm{\Omega }}^2(12\overline{m})=3\beta _0^2=3\overline{V}^2\mathrm{cn}^2(,\overline{m})`$ and $`\overline{V}^2=2\overline{m}\overline{\mathrm{\Omega }}^2`$, see Proof of Lemma 1.
(L.2) We observe that
$$\frac{d}{d\xi }\left[\frac{\mathrm{cn}(\xi )}{\mathrm{dn}(\xi )}\right]=\frac{(m1)\mathrm{sn}(\xi )}{\mathrm{dn}^2(\xi )},$$
then we can integrate by parts
$$_0^{4K}\frac{\mathrm{sn}^2(\xi )}{\mathrm{dn}^2(\xi )}๐\xi =_0^{4K}\frac{\mathrm{sn}(\xi )}{m1}\frac{d}{d\xi }\left[\frac{\mathrm{cn}(\xi )}{\mathrm{dn}(\xi )}\right]๐\xi =\frac{1}{1m}_0^{4K}\mathrm{cn}^2(\xi )๐\xi .$$
(L.3) We compute the derivative
$$\frac{d}{d\xi }\left[\frac{\mathrm{cn}(\xi )\mathrm{sn}(\xi )}{\mathrm{dn}(\xi )}\right]=2\mathrm{c}\mathrm{n}^2(\xi )1+m\frac{\mathrm{sn}^2(\xi )\mathrm{cn}^2(\xi )}{\mathrm{dn}^2(\xi )}$$
and integrate on the period $`[0,4K]`$.
(L.4) From the formula (68) of $`L`$ we have
$$\begin{array}{cc}\hfill gL[h]hL[g]=& \frac{d}{d\xi }\left[\left(_0^\xi h\overline{v}\right)\left(_0^\xi g\overline{u}\right)\right]\frac{d}{d\xi }\left[\left(_0^\xi h\overline{u}\right)\left(_0^\xi g\overline{v}\right)\right]+\hfill \\ & +\frac{1}{\overline{V}^2k}\mathrm{\hspace{0.17em}2}\pi \left[h\overline{v}g\overline{v}g\overline{v}h\overline{v}\right]=0.\hfill \end{array}$$
(L.5) By definition, $`L[1]`$ satisfies $`L[1]^{\prime \prime }+\left(3\beta _0^2+3\beta _0^2\right)L[1]=1`$, so we integrate on the period $`[0,2\pi ]`$.
(L.6),(L.7) Similarly by definition of $`L[\beta _0],L[\beta _0^2]`$; recall that $`\beta _0=0`$.
(L.8) By (L.6) and (L.4), it is sufficient to show that $`L[\beta _0]=0`$. From the formula (68), integrating by parts we have
$$L[\beta _0]=\beta _0\overline{v}\left(_0^\xi \overline{u}\right)\left(_0^\xi \beta _0\overline{u}\right)\overline{v}+\frac{1}{\overline{V}^2k}\left(_0^{2\pi }\beta _0\overline{v}\right)\overline{v}.$$
From the formulas (62), (64) of $`\overline{u},\overline{v},`$ recalling that $`\beta _0(\xi )=\overline{V}\mathrm{cn}(\overline{\mathrm{\Omega }}\xi )`$, we compute
$$_0^\xi \overline{u}=\frac{1}{\overline{\mathrm{\Omega }}^2}\left(\mathrm{cn}(\overline{\mathrm{\Omega }}\xi )1\right),_0^\xi \beta _0\overline{u}=\frac{\overline{V}}{2\overline{\mathrm{\Omega }}^2}\left(\mathrm{cn}^2(\overline{\mathrm{\Omega }}\xi )1\right).$$
(69)
Observe that $`_0^{2\pi }\mathrm{cn}(\overline{\mathrm{\Omega }}\xi )\frac{\mathrm{sn}^2(\overline{\mathrm{\Omega }}\xi )}{\mathrm{dn}^2(\overline{\mathrm{\Omega }}\xi )}๐\xi =0`$ by odd-symmetry with respect to $`\frac{\pi }{2}`$ on $`[0,\pi ]`$ and periodicity. So, recalling that $`\overline{V}^2=2\overline{m}\overline{\mathrm{\Omega }}^2`$, we compute $`\overline{v}=\frac{\overline{m}k}{\pi }`$. We can resume the computation of $`L[\beta _0]`$ obtaining
$$L[\beta _0]=\frac{3\overline{V}}{2\overline{\mathrm{\Omega }}^2}\overline{v}(\xi )\mathrm{cn}^2(\overline{\mathrm{\Omega }}\xi )\frac{\overline{V}\overline{m}k}{2\pi \overline{\mathrm{\Omega }}^2}.$$
Since $`\mathrm{cn}^3=\mathrm{cn}^3\frac{\mathrm{sn}^2}{\mathrm{dn}^2}=0`$ by the same odd-symmetry reason, by (64) we have $`\overline{v}(\xi )\mathrm{cn}^2(\overline{\mathrm{\Omega }}\xi )=\frac{\overline{m}k}{3\pi }`$, and so $`L[\beta _0]=0`$.
Moreover we can remark that by (L.4) there holds also $`\beta _0L[1]=0`$.
(L.9) By (L.5), it is equivalent to show that $`L[1]0`$. From the formula (68), integrating by parts we have
$$L[1]=\frac{2\pi }{\overline{V}^2k}\overline{v}^22\left(_0^\xi \overline{u}\right)\overline{v}.$$
We know that $`\overline{v}=\frac{\overline{m}k}{\pi }`$, so by (69)
$$L[1]=\frac{1}{\overline{\mathrm{\Omega }}^2}\overline{v}(\xi )\left(2\mathrm{c}\mathrm{n}(\overline{\mathrm{\Omega }}\xi )1\right).$$
ยฟFrom the equalities (L.1) and (L.3) we have $`\overline{v}(\xi )\mathrm{cn}(\overline{\mathrm{\Omega }}\xi )=\frac{2}{3}(12\overline{m})+\frac{\overline{m}k}{2\pi }`$, thus
$$L[1]=\frac{4(12\overline{m})}{3\overline{\mathrm{\Omega }}^2}$$
(70)
and this is strictly positive because $`\overline{m}<\frac{1}{2}`$.
(L.10) From (68) integrating by parts we have
$$\beta _0L[\beta _0]=2\beta _0\overline{v}\left(_0^\xi \beta _0\overline{u}\right)+\frac{2\pi }{\overline{V}^2k}\beta _0\overline{v}^2.$$
Using (L.3), integrating by parts and recalling the definition (66) of $`k`$ we compute
$$\beta _0\overline{v}=\overline{V}\overline{m}\mathrm{cn}^2+\frac{\overline{V}\overline{m}k}{2\pi }+\frac{\overline{V}(12\overline{m})}{2}$$
and, by (L.1) and (67),
$$\beta _0\overline{v}=\frac{\overline{V}(78\overline{m})}{12(1\overline{m})}.$$
(71)
By (69), $`\beta _0\overline{v}\left(_0^\xi \beta _0\overline{u}\right)=\frac{\overline{V}}{2\overline{\mathrm{\Omega }}^2}\beta _0\overline{v}\mathrm{cn}^2+\frac{\overline{V}}{2\overline{\mathrm{\Omega }}^2}\beta _0\overline{v}`$. The functions $`\beta _0`$ and $`\overline{v}`$ satisfy $`\beta _0^{\prime \prime }+\beta _0^3+3\beta _0^2\beta _0=0`$ and $`\overline{v}^{\prime \prime }+3\beta _0^2\overline{v}+3\beta _0^2\overline{v}=0`$, so that
$$\overline{v}^{\prime \prime }\beta _0\overline{v}\beta _0^{\prime \prime }+2\beta _0^3\overline{v}=0.$$
(72)
Deriving (65) we have $`\overline{v}^{}(2\pi )\overline{v}^{}(0)=\overline{V}^2k`$, so we can integrate (72) obtaining
$$\beta _0^3\overline{v}=\frac{\overline{V}^3k}{4\pi };$$
since $`\beta _0^3\overline{v}=\overline{V}^2\beta _0\overline{v}\mathrm{cn}^2`$, we write
$$\beta _0\overline{v}\left(_0^\xi \beta _0\overline{u}\right)=\frac{\overline{m}k}{4\pi }+\frac{\overline{V}}{2\overline{\mathrm{\Omega }}^2}\beta _0\overline{v}.$$
Thus, by (71) and (66), we can express $`\beta _0L[\beta _0]`$ in terms of $`\overline{m}`$ only,
$$\beta _0L[\beta _0]=\frac{32\overline{m}^232\overline{m}1}{12(16\overline{m}^216\overline{m}+1)}=\frac{1}{6}\frac{1}{4(16\overline{m}^216\overline{m}+1)}.$$
(73)
The polynomial $`p(m)=16m^216m+1`$ is non-zero for $`m(\frac{2\sqrt{3}}{4},\frac{2+\sqrt{3}}{4})`$ and $`\overline{m}(0.20,0.21)`$; so $`B_0=\frac{6}{4p(\overline{m})}0`$, in particular $`B_0(1,0.9)`$.
(L.11) From (73) it follows that $`C_00`$, in particular $`2.9<C_0=2\frac{3}{2p(\overline{m})}<3`$.
(L.12) By Exchange rule (L.4) and (L.5), it is sufficient to show that $`A_01`$, that is $`3\beta _0^2L[1]1`$. Recall that, by construction of $`\overline{m}`$, $`3\beta _0^2=\overline{\mathrm{\Omega }}^2(12\overline{m})`$. So from (70) it follows
$$3\beta _0^2L[1]=\frac{4}{3}(12\overline{m})^2,$$
and $`\frac{4}{3}(12m)^2=1`$ if and only if $`16m^216m+1=0`$, while $`\overline{m}(0.20,0.21)`$, like above; in particular $`0.4<\mathrm{\hspace{0.17em}3}\beta _0^2L[1]<\mathrm{\hspace{0.17em}0.5}.`$ $`\mathrm{}`$
*Remark.* Approximated computations give
$$\begin{array}{ccc}\overline{m}(0.20,\mathrm{\hspace{0.17em}0.21})\hfill & \overline{\sigma }(2.10,\mathrm{\hspace{0.17em}2.16})\hfill & \overline{\mathrm{\Omega }}(1.05,\mathrm{\hspace{0.17em}1.06})\hfill \\ \overline{V}^2(0.44,\mathrm{\hspace{0.17em}0.48})\hfill & \mathrm{cn}^2(2.85,\mathrm{\hspace{0.17em}2.90})\hfill & \beta _0^2(1.27,\mathrm{\hspace{0.17em}1.37}).\hfill \end{array}$$
Lemma 4. *The partial derivative $`_ZG(1,0,\beta _0,\beta _0)`$ is an invertible operator.*
*Proof.* Let $`_ZG(1,0,\beta _0,\beta _0)[\eta ,h,k]=(0,0,0)`$ for some $`(\eta ,h,k)Z,`$ that is
$`\begin{array}{c}6\eta \beta _0^2+3\beta _0^2h+3\beta _0^2k=0\hfill \\ 3\eta \left(\beta _0^2\beta _0^2\right)+h^{\prime \prime }+\left(3\beta _0^2+3\beta _0^2\right)h3\beta _0^2h+6\beta _0k\beta _0=0\hfill \\ 3\eta \left(\beta _0^2\beta _0^2\right)+k^{\prime \prime }+\left(3\beta _0^2+3\beta _0^2\right)k3\beta _0^2k+6\beta _0h\beta _0=0.\hfill \end{array}`$ (77)
We evaluate the second and the third equation at the same variable and subtract; $`\rho =hk`$ satisfies
$`\rho ^{\prime \prime }+\left(3\beta _0^2+3\beta _0^2\right)\rho 3\beta _0^2\rho 6\beta _0\rho \beta _0=0.`$ (78)
By definition of $`L`$, see Lemma 2, (78) can be written as
$`\rho =3\beta _0^2\rho L[1]+6\beta _0\rho L[\beta _0].`$ (79)
Multiplying this equation by $`\beta _0^2`$ and integrating we obtain
$`\beta _0^2\rho \left(13\beta _0^2L[1]\right)=6\beta _0\rho \beta _0^2L[\beta _0].`$
In Lemma 3 we prove that $`\left(13\beta _0^2L[1]\right)=A_00`$ and $`\beta _0^2L[\beta _0]=0,`$ then $`\beta _0^2\rho =0`$.
On the other hand, multiplying (79) by $`\beta _0`$ and integrating we have
$`\beta _0\rho \left(16\beta _0L[\beta _0]\right)=3\beta _0^2\rho \beta _0L[1];`$
in Lemma 3 we show that $`\left(16\beta _0L[\beta _0]\right)=B_00`$ and $`\beta _0L[1]=0,`$ then $`\beta _0\rho =0.`$ From (79) we have so $`\rho =0.`$ Thus $`h=k`$ and (77) becomes
$`\begin{array}{c}\eta \beta _0^2+\beta _0^2h=0\hfill \\ 3\eta \left(\beta _0^2\beta _0^2\right)+h^{\prime \prime }+\left(3\beta _0^2+3\beta _0^2\right)h3\beta _0^2h+6\beta _0h\beta _0=0.\hfill \end{array}`$
By substitution we have
$$h=3\eta L[\beta _0^2]6\beta _0hL[\beta _0].$$
Multiplying, as before, by $`\beta _0^2`$ and by $`\beta _0`$ and integrating, we obtain $`\beta _0h=\beta _0^2h=0`$ because $`\left(1+6\beta _0L[\beta _0]\right)=C_00,`$$`\beta _0L[\beta _0^2]=0,`$ and $`\beta _0^23\beta _0^2L[\beta _0^2]=3\beta _0^2L[\beta _0^2]0`$, see Lemma 3 again. Thus $`h=0`$, $`\eta =0`$ and the derivative $`_ZG(1,0,\beta _0,\beta _0)`$ is injective.
The operator $`ZZ,(\eta ,h,k)((6\beta _0^2)^1\eta ,L[h],L[k])`$ is compact because $`L`$ is compact, see Lemma 2. So, by the Fredholm Alternative, the partial derivative $`_ZG(1,0,\beta _0,\beta _0)`$ is also surjective. $`\mathrm{}`$
By the Implicit Function Theorem and the regularity of $`G`$, using the rescaling (52) we obtain, for $`|b\frac{1}{2}|`$ and $`\epsilon `$ small enough, the existence of a solution close to $`(0,\beta _0,\beta _0)`$ for the $`Z`$-equation (33).
More precisely: from Lemma 1 and 4 it follows the existence of a $`๐^1`$-function $`g`$ defined on a neighborhood of $`\lambda =1`$ such that
$$G(\lambda ,g(\lambda ))=0,$$
that is, $`g(\lambda )`$ solves (56), and $`g(1)=(0,\beta _0,\beta _0)`$. Moreover, for $`|\lambda 1|`$ small, it holds
$$g(\lambda )g(1)_\sigma \stackrel{~}{c}|\lambda 1|$$
(81)
for some positive constant $`\stackrel{~}{c}`$. In the following, we denote several positive constants with the same symbol $`\stackrel{~}{c}`$.
We set $`\mathrm{\Phi }_{(b,\epsilon )}:(\widehat{u}_{0,0},r,s)(c,x,y)`$ the rescaling map (52) and $`H_{(b,\epsilon )}:\times ZZ`$ the operator corresponding to the auxiliary bifurcation equation (45), which so can be written as
$$H_{(b,\epsilon )}(\mu ,z)=0.$$
We define
$$z_{(b,\epsilon )}^{}=\mathrm{\Phi }_{(b,\epsilon )}^1[g(\lambda _{(b,\epsilon )})],$$
thus it holds $`H_{(b,\epsilon )}(0,z_{(b,\epsilon )}^{})=0`$, that is, $`z_{(b,\epsilon )}^{}`$ solves the bifurcation equation (45) for $`\mu =0`$.
We observe that $`p_{(b,\epsilon )}(0,z)=_zp_{(b,\epsilon )}(0,z)=0`$ for every $`z`$ and so, in particular, for $`z=z_{(b,\epsilon )}^{}`$; it follows that
$$_zH_{(b,\epsilon )}(0,z_{(b,\epsilon )}^{})=(\mathrm{\Phi }_{(b,\epsilon )}^1)^3_zG(\lambda _{(b,\epsilon )},g(\lambda _{(b,\epsilon )}))\mathrm{\Phi }_{(b,\epsilon )}.$$
(82)
$`G`$ is of class $`๐^1`$, so $`_ZG(\lambda ,g(\lambda ))`$ remains invertible for $`\lambda `$ sufficiently close to 1. Notice that $`\lambda _{(b,\epsilon )}`$ is sufficiently close to 1 if $`|b\frac{1}{2}|`$ and $`\epsilon `$ are small enough. Then, by (82), the partial derivative $`_zH_{(b,\epsilon )}(0,z_{(b,\epsilon )}^{})`$ is invertible. By the Implicit Function Theorem, it follows that for every $`\mu `$ sufficiently small there exists a solution $`z_{(b,\epsilon )}(\mu )`$ of equation (45), that is
$$H_{(b,\epsilon )}(\mu ,z_{(b,\epsilon )}(\mu ))=0.$$
We indicate $`z_0=(0,\beta _0,\beta _0)`$. The operators $`\left(_zH_{(b,\epsilon )}(\mu ,z)\right)^1`$ and $`_\mu H_{(b,\epsilon )}(\mu ,z)`$ are bounded by some constant for every $`(\mu ,z)`$ in a neighborhood of $`(0,z_0),`$ uniformly in $`(b,\epsilon )`$, if $`|b1/2|,\epsilon `$ are small enough. So the implicit functions $`z_{(b,\epsilon )}`$ are defined on some common interval $`(\mu _0,\mu _0)`$ for $`|b1/2|`$, $`\epsilon `$ small, and it holds
$$z_{(b,\epsilon )}(\mu )z_{(b,\epsilon )}^{}_\sigma \stackrel{~}{c}|\mu |$$
(83)
for some $`\stackrel{~}{c}`$ which does not depend on $`(b,\epsilon )`$.
Such a common interval $`(\mu _0,\mu _0)`$ permits the evaluation $`z_{(b,\epsilon )}(\mu )`$ at $`\mu =\epsilon `$ for $`\epsilon <\mu _0`$, obtaining a solution of the original bifurcation equation written in (33).
Moreover, $`\mathrm{\Phi }_{(b,\epsilon )}^1\mathrm{Id}_Z_\sigma =|\sqrt{b(2+b\epsilon ^2)}1||b\frac{1}{2}|+\epsilon ^2`$, so, by (81) and triangular inequality,
$`\begin{array}{c}z_{(b,\epsilon )}^{}z_0_\sigma \stackrel{~}{c}(|b\frac{1}{2}|+\epsilon ^2).\end{array}`$ (85)
Thus from (83) and (85) we have
$`\begin{array}{c}z_{(b,\epsilon )}(\epsilon )z_0_\sigma \stackrel{~}{c}(|b\frac{1}{2}|+\epsilon ),\end{array}`$
and, by (38),
$`\begin{array}{c}p(\epsilon ,z_{(b,\epsilon )}(\epsilon ))_\sigma \stackrel{~}{c}\epsilon .\end{array}`$
*Remark.* Since the solutions $`z_{(b,\epsilon )}(\epsilon )`$ are close to $`z_0=(0,\beta _0,\beta _0)`$, they actually depend on the two arguments $`(\phi _1,\phi _2)`$; this is a necessary condition for the quasi-periodicity.
We define $`u_{(b,\epsilon )}=z_{(b,\epsilon )}(\epsilon )+p_{(b,\epsilon )}(\epsilon ,z_{(b,\epsilon )}(\epsilon ))`$. Renaming $`\mu _0=\epsilon _0`$, we have finally proved:
Theorem 1. *Let $`\overline{\sigma }>0`$, $`\beta _0`$ as in Lemma 1, $`\stackrel{~}{B}_\gamma `$ as in (34) with $`\gamma (0,\frac{1}{4})`$.* *For every $`\sigma (0,\overline{\sigma })`$, there exist positive constants $`\delta _0`$, $`\epsilon _0`$, $`\overline{c}_1`$, $`\overline{c}_2`$ and the uncountable Cantor set*
$$_\gamma =\{(b,\epsilon )(\frac{1}{2}\delta _0,\frac{1}{2}+\delta _0)\times (0,\epsilon _0):\frac{\epsilon }{2+b\epsilon ^2},b\epsilon ^2\stackrel{~}{B}_\gamma ,\frac{1+b\epsilon ^2}{\epsilon }\}$$
*such that, for every $`(b,\epsilon )_\gamma `$, there exists a solution $`u_{(b,\epsilon )}_\sigma `$ of (13).* *According to decomposition (23), $`u_{(b,\epsilon )}`$ can be written as*
$$u_{(b,\epsilon )}(\phi _1,\phi _2)=\widehat{u}_{0,0}+r(\phi _1)+s(\phi _2)+p(\phi _1,\phi _2),$$
*where its components satisfy*
$$\begin{array}{c}r\beta _0_\sigma +s\beta _0_\sigma +|\widehat{u}_{0,0}|\overline{c}_1(|b\frac{1}{2}|+\epsilon ),p_\sigma \overline{c}_2\epsilon .\end{array}$$
*As a consequence, problem (3) admits uncountable many small amplitude, analytic, quasi-periodic solutions $`v_{(b,\epsilon )}`$ with two frequencies, of the form (7):*
$$\begin{array}{cc}\hfill v_{(b,\epsilon )}(t,x)& =\epsilon u_{(b,\epsilon )}(\epsilon t,(1+b\epsilon ^2)t+x)\hfill \\ & =\epsilon \left[\widehat{u}_{0,0}+r(\epsilon t)+s((1+b\epsilon ^2)t+x)+๐ช(\epsilon )\right]\hfill \\ & =\epsilon \left[\beta _0(\epsilon t)+\beta _0((1+b\epsilon ^2)t+x)+๐ช\left(|b\frac{1}{2}|+\epsilon \right)\right].\hfill \end{array}$$
## 6. Waves traveling in opposite directions
In this section we look for solutions of (3) of the form (4),
$$v(t,x)=u(\omega _1t+x,\omega _2tx),$$
for $`u_\sigma `$. We introduce two parameters $`(a,\epsilon )^2`$ and set the frequencies as in ,
$$\omega _1=1+\epsilon ,\omega _2=1+a\epsilon .$$
For functions of the form (4), problem (3) is written as
$$L_{a,\epsilon }[u]=u^3+f(u)$$
where
$$L_{a,\epsilon }=\epsilon (2+\epsilon )_{\phi _1}^2+2\left(2+(a+1)\epsilon +a\epsilon ^2\right)_{\phi _1\phi _2}^2+a\epsilon (2+a\epsilon )_{\phi _2}^2.$$
We rescale $`u\sqrt{\epsilon }u`$ and define $`f_\epsilon (u)=\epsilon ^{3/2}f(\sqrt{\epsilon }u)`$. Thus the problem can be written as
$$L_{a,\epsilon }[u]=\epsilon u^3+\epsilon f_\epsilon (u).$$
(88)
For $`\epsilon =0`$, the operator is $`L_{a,0}=4_{\phi _1\phi _2}^2`$; its kernel is the direct sum $`Z=CQ_1Q_2`$, see (20). Writing $`u`$ in Fourier series we obtain an expression similar to (14),
$`L_{a,\epsilon }[u]={\displaystyle \underset{(m,n)^2}{}}D_{a,\epsilon }(m,n)\widehat{u}_{mn}e^{im\phi _1}e^{in\phi _2},`$
where the eigenvalues $`D_{a,\epsilon }(m,n)`$ are given by
$`D_{a,\epsilon }(m,n)=`$ $`\epsilon (2+\epsilon )m^2+a\epsilon (2+a\epsilon )n^2+2\left(2+(a+1)\epsilon +a\epsilon ^2\right)mn`$
$`=`$ $`(2+\epsilon )(2+a\epsilon )\left(m+{\displaystyle \frac{a\epsilon }{2+\epsilon }}n\right)\left({\displaystyle \frac{\epsilon }{2+a\epsilon }}m+n\right).`$
By Lyapunov-Schmidt reduction we project the equation (88) on the four subspaces,
$`\begin{array}{cc}\hfill 0=& \widehat{u}_{0,0}^3+3\widehat{u}_{0,0}\left(r^2+s^2\right)+r^3+s^3+\hfill \\ & +\mathrm{\Pi }_C\left[(u^3z^3)f_\epsilon (u)\right]\left[Cequation\right]\hfill \\ \hfill (2+\epsilon )r^{\prime \prime }=& 3\widehat{u}_{0,0}^2r+3\widehat{u}_{0,0}\left(r^2r^2\right)+r^3r^3+3s^2r+\hfill \\ & +\mathrm{\Pi }_{Q_1}\left[(u^3z^3)f_\epsilon (u)\right]\left[Q_1equation\right]\hfill \\ \hfill a(2+a\epsilon )s^{\prime \prime }=& 3\widehat{u}_{0,0}^2s+3\widehat{u}_{0,0}\left(s^2s^2\right)+s^3s^3+3r^2s+\hfill \\ & +\mathrm{\Pi }_{Q_2}\left[(u^3z^3)f_\epsilon (u)\right]\left[Q_2equation\right]\hfill \\ \hfill L_{a,\epsilon }[p]=& \epsilon \mathrm{\Pi }_P\left[u^3+f_\epsilon (u)\right].\left[Pequation\right]\hfill \end{array}`$
We repeat the arguments of Appendix C and find a Cantor set $`๐_\gamma `$ such that $`|D_{a,\epsilon }(m,n)|>\gamma `$ for every $`(a,\epsilon )๐_\gamma `$. Then $`L_{a,\epsilon }`$ is invertible for $`(a,\epsilon )๐_\gamma `$ and the $`P`$-equation can be solved as in the section 4.
We repeat the same procedure already shown in section 5 and solve the bifurcation equation. The only differences are:
\- the parameter $`a`$ tends to $`1`$ instead of $`b\frac{1}{2}`$;
\- the rescaling map is $`\mathrm{\Psi }_{(a,\epsilon )}:(\widehat{u}_{0,0},r,s)(c,x,y)`$, where
$`\begin{array}{cc}r=\sqrt{2+\epsilon }x\hfill & \widehat{u}_{0,0}=\sqrt{2+\epsilon }c\hfill \\ s=\sqrt{a(2+a\epsilon )}y\hfill & \lambda =\lambda _{(a,\epsilon )}=\frac{a(2+a\epsilon )}{2+\epsilon },\hfill \end{array}`$
instead of $`\mathrm{\Phi }_{(b,\epsilon )}`$ defined in (52).
We note that by means of the rescaling map $`\mathrm{\Psi }_{(a,\epsilon )}`$ we obtain just the equation (56). Thus we conclude:
Theorem 2. *Let $`\overline{\sigma }>0`$, $`\beta _0`$ as in Lemma 1, $`\stackrel{~}{B}_\gamma `$ as in (34) with $`\gamma (0,\frac{1}{4})`$.* *For every $`\sigma (0,\overline{\sigma })`$, there exist positive constants $`\delta _0`$, $`\epsilon _0`$, $`\overline{c}_1`$, $`\overline{c}_2`$ and the uncountable Cantor set*
$$๐_\gamma =\{(a,\epsilon )(1\delta _0,\mathrm{\hspace{0.17em}1}+\delta _0)\times (0,\epsilon _0):\frac{a\epsilon }{2+\epsilon },\frac{\epsilon }{2+a\epsilon }\stackrel{~}{B}_\gamma ,\frac{1+\epsilon }{1+a\epsilon }\}$$
*such that, for every $`(a,\epsilon )๐_\gamma `$, there exists a solution $`u_{(a,\epsilon )}_\sigma `$ of (88).* *According to decomposition (23), $`u_{(a,\epsilon )}`$ can be written as*
$$u_{(a,\epsilon )}(\phi _1,\phi _2)=\widehat{u}_{0,0}+r(\phi _1)+s(\phi _2)+p(\phi _1,\phi _2),$$
*where its components satisfy*
$$r\beta _0_\sigma +s\beta _0_\sigma +|\widehat{u}_{0,0}|\overline{c}_1(|a1|+\epsilon ),p_\sigma \overline{c}_2\epsilon .$$
*As a consequence, problem (3) admits uncountable many small amplitude, analytic, quasi-periodic solutions $`v_{(a,\epsilon )}`$ with two frequencies, of the form (4):*
$$\begin{array}{cc}\hfill v_{(a,\epsilon )}(t,x)& =\sqrt{\epsilon }u_{(a,\epsilon )}((1+\epsilon )t+x,(1+a\epsilon )tx)\hfill \\ & =\sqrt{\epsilon }\left[\widehat{u}_{0,0}+r\left((1+\epsilon )t+x\right)+s\left((1+a\epsilon )tx\right)+๐ช(\epsilon )\right]\hfill \\ & =\sqrt{\epsilon }\left[\beta _0((1+\epsilon )t+x)\right)+\beta _0((1+a\epsilon )tx)+๐ช(|a1|+\epsilon )].\hfill \end{array}$$
## 7. Appendix A. Hilbert algebra property of $`_\sigma `$
Let $`u,v_\sigma `$, $`u=_{m^2}\widehat{u}_me^{im\phi }`$, $`v=_{m^2}\widehat{v}_me^{im\phi }`$. The product $`uv`$ is
$$uv=\underset{j}{}\left(\underset{k}{}\widehat{u}_{jk}\widehat{v}_k\right)e^{ij\phi },$$
so its $`_\sigma `$-norm, if it converges, is
$$uv_\sigma ^2=\underset{j}{}\left|\underset{k}{}\widehat{u}_{jk}\widehat{v}_k\right|^2\left(1+|j|^{2s}\right)e^{2|j|\sigma }.$$
We define
$$a_{jk}=\left[\frac{\left(1+|jk|^{2s}\right)\left(1+|k|^{2s}\right)}{\left(1+|j|^{2s}\right)}\right]^{1/2}.$$
Given any $`(x_k)_k`$, it holds by Hรถlder inequality
$$\left|\underset{k}{}x_k\right|^2=\left|\underset{k}{}\frac{1}{a_{jk}}x_ka_{jk}\right|^2c_j^2\underset{k}{}|x_ka_{jk}|^2,$$
(91)
where
$$c_j^2:=\underset{k}{}\frac{1}{a_{jk}^2}.$$
We show that there exists a constant $`c>0`$ such that $`c_j^2c^2`$ for every $`j^2`$. We recall that, fixed $`p1`$, it holds
$$(a+b)^p2^{p1}\left(a^p+b^p\right)a,b0.$$
Then, for $`s\frac{1}{2}`$, we have
$$\begin{array}{c}1+|j|^{2s}1+\left(|jk|+|k|\right)^{2s}1+2^{2s1}\left(|jk|^{2s}+|k|^{2s}\right)\hfill \\ \hfill <2^{2s1}\left(1+|jk|^{2s}+1+|k|^{2s}\right),\end{array}$$
so
$$\frac{1}{a_{jk}^2}<2^{2s1}\left(\frac{1}{1+|jk|^{2s}}+\frac{1}{1+|k|^{2s}}\right).$$
The series $`_{k^2}\frac{1}{1+|k|^p}`$ converges if $`p>2`$, thus for $`s>1`$
$$\begin{array}{c}c_j^2<2^{2s1}\left(\underset{k}{}\frac{1}{1+|jk|^{2s}}+\underset{k}{}\frac{1}{1+|k|^{2s}}\right)=2^{2s}\underset{k^2}{}\frac{1}{1+|k|^{2s}}:=c^2<\mathrm{}.\hfill \end{array}$$
We put $`x_k=\widehat{u}_{jk}\widehat{v}_k`$ in (91) and compute
$`\left|{\displaystyle \underset{k}{}}\widehat{u}_{jk}\widehat{v}_k\right|^2\left(1+|j|^{2s}\right)`$ $`c^2{\displaystyle \underset{k}{}}|\widehat{u}_{jk}\widehat{v}_ka_{jk}|^2\left(1+|j|^{2s}\right)`$
$`=c^2{\displaystyle \underset{k}{}}|\widehat{u}_{jk}\widehat{v}_k|^2\left(1+|jk|^{2s}\right)\left(1+|k|^{2s}\right),`$
$`uv_\sigma ^2`$ $`={\displaystyle \underset{j}{}}\left|{\displaystyle \underset{k}{}}\widehat{u}_{jk}\widehat{v}_k\right|^2\left(1+|j|^{2s}\right)e^{2|j|\sigma }`$
$`{\displaystyle \underset{j}{}}c^2{\displaystyle \underset{k}{}}|\widehat{u}_{jk}|^2|\widehat{v}_k|^2\left(1+|jk|^{2s}\right)\left(1+|k|^{2s}\right)e^{2(|jk|+|k|)\sigma }`$
$`=c^2{\displaystyle \underset{k}{}}\left({\displaystyle \underset{j}{}}|\widehat{u}_{jk}|^2\left(1+|jk|^{2s}\right)e^{2|jk|\sigma }\right)|\widehat{v}_k|^2\left(1+|k|^{2s}\right)e^{2|k|\sigma }`$
$`=c^2u_\sigma ^2v_\sigma ^2.`$
So $`uv_\sigma cu_\sigma v_\sigma `$ for all $`u,v_\sigma `$. We notice that the constant $`c`$ depends only on $`s`$,
$$c=2^s\left(\underset{k^2}{}\frac{1}{1+|k|^{2s}}\right)^{1/2}.$$
## 8. Appendix B. Change of form for quasi-periodic functions
First an algebric proposition, then we show that one can pass from (5) to (7) without loss of generality.
Proposition. *Let $`A,B\mathrm{Mat}_2()`$ be invertible matrices such that $`AB^1`$ has integer coefficient. Then, given any $`u_\sigma `$, the function $`v(t,x)=u\left(A(t,x)\right)`$ can be written as $`v(t,x)=w\left(B(t,x)\right)`$ for some $`w_\sigma `$, that is*$`\{uA:u_\sigma \}\{wB:w_\sigma \}`$.
*Proof.* Let $`u_\sigma `$. The function $`uA`$ belongs to $`\{wB:w_\sigma \}`$ if and only if $`uAB^1=w`$ for some $`w_\sigma `$, and this is true if and only if $`uAB^1`$ is $`2\pi `$ periodic; since $`AB^1\mathrm{Mat}_2()`$, we can conclude. $`\mathrm{}`$
Lemma. *The set of the quasi-periodic functions of the form (5) is equal to the set of the quasi-periodic functions of the form (7), that is,*
$`\begin{array}{c}\{v:v(t,x)=u(\omega _1t+x,\omega _2t+x),(\omega _1,\omega _2)^2,\omega _10,\omega _20,\frac{\omega _1}{\omega _2},u_\sigma \}\hfill \\ =\{v:v(t,x)=u(\epsilon t,(1+b\epsilon ^2)t+x),(b,\epsilon )^2,\epsilon 0,(1+b\epsilon ^2)0,\hfill \\ \frac{1+b\epsilon ^2}{\epsilon },u_\sigma \}.\hfill \end{array}`$
*Proof.* Given any $`\omega _1,\omega _2,b,\epsilon `$, we define
$`A=\left(\begin{array}{cc}\omega _1& 1\\ \omega _2& 1\end{array}\right)B=\left(\begin{array}{cc}\epsilon & 0\\ (1+b\epsilon ^2)& 1\end{array}\right).`$
Let $`v(t,x)`$ be any element of the set of quasi-periodic functions of the form (5), that is $`v=uA`$ for some fixed $`\omega _1,\omega _20`$ such that $`\frac{\omega _1}{\omega _2}`$ and $`u_\sigma `$. We observe that $`v`$ belongs to the set of quasi-periodic functions of the form (7) if $`v=wB`$ for some $`(b,\epsilon )`$ such that $`\epsilon 0,\frac{1+b\epsilon ^2}{\epsilon }`$ and some $`w_\sigma `$. By the Proposition, this is true if we find $`(b,\epsilon )`$ such that $`AB^1\mathrm{Mat}_2()`$. We can choose
$$b=\frac{\omega _21}{(\omega _1\omega _2)^2},\epsilon =\omega _1\omega _2,$$
so that $`AB^1=\left(\begin{array}{cc}1& 1\\ 0& 1\end{array}\right)`$. We notice that $`\frac{1+b\epsilon ^2}{\epsilon }`$ if and only if $`\frac{\omega _1}{\omega _2}`$.
Conversely, we fix $`(b,\epsilon )`$ and look for $`(\omega _1,\omega _2)`$ such that $`BA^1\mathrm{Mat}_2()`$. This condition is satisfied if we choose the inverse transformation, $`\omega _1=1+\epsilon +b\epsilon ^2`$, $`\omega _2=1+b\epsilon ^2`$. $`\mathrm{}`$
## 9. Appendix C. Small divisors
Fixed $`\gamma (0,\frac{1}{4})`$, we have defined in (34) the set $`\stackrel{~}{B}_\gamma `$ of โbadly approximable numbersโ as
$`\stackrel{~}{B}_\gamma =\{x:\left|m+nx\right|>{\displaystyle \frac{\gamma }{|n|}}m,n,m0,n0\}.`$
$`\stackrel{~}{B}_\gamma `$ is non-empty, symmetric, it has zero Lebesgue-measure and it accumulates to 0. Moreover, for every $`\delta >0`$, $`\stackrel{~}{B}_\gamma (\delta ,\delta )`$ is uncountable.
In fact, $`\stackrel{~}{B}_\gamma `$ contains all the irrational numbers whose continued fractions expansion is of the form $`[0,a_1,a_2,\mathrm{}]`$, with $`a_j<\gamma ^12`$ for every $`j2`$. Such a set is uncountable: since $`\gamma ^12>2`$, for every $`j1`$ there are at least two choices for the value of $`a_j`$. Moreover, it accumulates to 0: if $`y=[0,a_1,a_2,\mathrm{}]`$, it holds $`0<y<a_1^1`$, and $`a_1`$ has no upper bound. See also Remark 2.4 in and, for the inclusion of such a set in $`\stackrel{~}{B}_\gamma `$, the proof of Theorem 5F in , p. 22.
We prove the following estimate.
Proposition. *Let $`\gamma (0,\frac{1}{4}),\delta (0,\frac{1}{2}).`$ Then for all $`x,y\stackrel{~}{B}_\gamma (\delta ,\delta )`$ it holds*
$`\left|(m+nx)(my+n)\right|>\gamma (1\delta \delta ^2)m,n,m,n0.`$
*Proof.* We shortly set $`D=\left|(m+nx)(my+n)\right|.`$ There are four cases.
Case 1. $`|m+nx|>1,|my+n|>1.`$ Then $`|D|>1`$.
Case 2. $`|m+nx|<1,|my+n|>1.`$ Multiplying the first inequality by $`|y|,`$
$`\begin{array}{cc}\hfill |y|>& |my+nxy|=|my+nn(1xy)|\hfill \\ \hfill & \left||n(1xy)||my+n|\right||n(1xy)||my+n|,\hfill \end{array}`$
so $`|my+n|>|n|(1xy)|y|`$ and
$`\begin{array}{cc}\hfill |D|>& \frac{\gamma }{|n|}\left[|n|(1xy)|y|\right]=\gamma \left[(1xy)\frac{|y|}{|n|}\right]\hfill \\ \hfill >& \gamma [(1\delta ^2)\delta ].\hfill \end{array}`$
Case 3. $`|m+nx|>1,|my+n|<1.`$ Analogous to case 2.
Case 4. $`|m+nx|<1,|my+n|<1.`$ Dividing the first inequality by $`|n|,`$ for triangular inequality we have
$$\left|\frac{m}{n}\right|\left|\frac{m}{n}+x\right|+|x|<\frac{1}{|n|}+\delta ,$$
and similarly $`\left|\frac{n}{m}\right|<\frac{1}{|m|}+\delta `$. So
$$\left(\frac{1}{|n|}+\delta \right)\left(\frac{1}{|m|}+\delta \right)>\left|\frac{n}{m}\frac{m}{n}\right|=1.$$
If $`|n|,|m|2,`$ then $`\left(\frac{1}{|n|}+\delta \right)\left(\frac{1}{|m|}+\delta \right)<1,`$ a contradiction. It follows that at least one between $`|n|`$ and $`|m|`$ is equal to 1. Suppose $`|n|=1`$. Then $`|m+nx|=|m\pm x||m|\delta `$ and
$$|D|>\frac{\gamma }{|m|}\left(|m|\delta \right)=\gamma \left(1\frac{\delta }{|m|}\right)\gamma (1\delta ).$$
If $`|m|=1`$ the conclusion is the same. $`\mathrm{}`$
Fixed $`\gamma (0,\frac{1}{4})`$ and $`\delta (0,\frac{1}{2})`$, we define the set
$`\begin{array}{cc}\hfill B(\gamma ,\delta )=\{(b,\epsilon )^2:& \epsilon 0,1+b\epsilon ^20,2+b\epsilon ^20,\hfill \\ & \frac{1+b\epsilon ^2}{\epsilon },\frac{\epsilon }{2+b\epsilon ^2},b\epsilon ^2\stackrel{~}{B}_\gamma (\delta ,\delta )\}\hfill \end{array}`$
and the map
$$g:B(\gamma ,\delta )^2,g(b,\epsilon )=(\frac{\epsilon }{2+b\epsilon ^2},b\epsilon ^2).$$
$`g^{}(b,\epsilon )`$ is invertible on $`B(\gamma ,\delta ).`$ Its imagine is the set $`R(g)=\{(x,y)^2:x,y\stackrel{~}{B}_\gamma (\delta ,\delta ),\frac{1}{x}y\}`$ and its inverse is
$`g^1(x,y)=({\displaystyle \frac{y(1xy)}{2x}},{\displaystyle \frac{2x}{1xy}}).`$
Thus $`B(\gamma ,\delta )`$ is homeomorphic to $`R(g)=\{(x,y)\stackrel{~}{B}_\gamma ^2:|x|,|y|<\delta ,\frac{1}{x}y\}.`$ We observe that, fixed any $`\overline{x}\stackrel{~}{B}_\gamma (\delta ,\delta )`$, it occurs $`\frac{1}{\overline{x}}y`$ only for countably many numbers $`y`$. We know that $`\stackrel{~}{B}_\gamma (\delta ,\delta )`$ is uncountable so, removing from $`[\stackrel{~}{B}_\gamma (\delta ,\delta )]^2`$ the couples $`\{(\overline{x},y):y=\frac{1}{\overline{x}}qq\}`$, it remains uncountably many other couples. Thus $`R(g)`$ is uncountable and so, through $`g`$, also $`B(\gamma ,\delta )`$.
Moreover, if we consider couples $`(x,y)[\stackrel{~}{B}_\gamma (\delta ,\delta )]^2`$ such that $`x0`$ and $`(x/y)1,`$ applying $`g^1`$ we find couples $`(b,\epsilon )B(\gamma ,\delta )`$ which satisfy $`\epsilon 0`$, $`b1/2`$. In other words, the set $`B(\gamma ,\delta )`$ accumulates to $`(1/2,0)`$.
Finally we estimate $`D_{b,\epsilon }(m,n)`$ for $`(b,\epsilon )B(\gamma ,\delta )`$. We have
$$|2+b\epsilon ^2|=\frac{2}{|1xy|}>\frac{2}{1+\delta ^2},$$
so from the previous Proposition and (3) it follows
$$\left|D_{b,\epsilon }(m,n)\right|=|D||2+b\epsilon ^2|>\gamma (1\delta \delta ^2)\frac{2}{1+\delta ^2}.$$
The factor on the right of $`\gamma `$ is greater than 1 if we choose, for example, $`\delta =1/4;`$ we define $`_\gamma =B(\gamma ,\delta )_{|\delta =\frac{1}{4}}`$, so that there holds
$$\left|D_{b,\epsilon }(m,n)\right|>\gamma (b,\epsilon )_\gamma .$$
We can observe that the condition $`\frac{1+b\epsilon ^2}{\epsilon }`$ implies $`1+b\epsilon ^20`$, that $`\frac{\epsilon }{2+b\epsilon ^2}\stackrel{~}{B}_\gamma `$ implies $`\epsilon 0`$ and $`|b\epsilon ^2|<\delta `$ implies $`2+b\epsilon ^20`$, so that we can write
$$_\gamma =\{(b,\epsilon )^2:\frac{\epsilon }{2+b\epsilon ^2},b\epsilon ^2\stackrel{~}{B}_\gamma ,\left|\frac{\epsilon }{2+b\epsilon ^2}\right|,|b\epsilon ^2|<\frac{1}{4},\frac{1+b\epsilon ^2}{\epsilon }\}.$$
We notice also that, for $`|b\frac{1}{2}|`$ and $`\epsilon `$ small enough, there holds automatically $`|\frac{\epsilon }{2+b\epsilon ^2}|<\frac{1}{4}`$, $`|b\epsilon ^2|<\frac{1}{4}`$. So, if we are interested to couples $`(b,\epsilon )`$ close to $`(\frac{1}{2},0)`$, say $`|b\frac{1}{2}|<\delta _0`$, $`|\epsilon |<\epsilon _0`$, we can write
$$_\gamma =\{(b,\epsilon )(\frac{1}{2}\delta _0,\frac{1}{2}+\delta _0)\times (0,\epsilon _0):\frac{\epsilon }{2+b\epsilon ^2},b\epsilon ^2\stackrel{~}{B}_\gamma ,\frac{1+b\epsilon ^2}{\epsilon }\}.$$
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# Polarization microvariability of BL Lac objects Based on observations made at the Complejo Astronรณmico El Leoncito, which is operated under agreement between CONICET and the National Universities of La Plata, Cรณrdoba, and San Juan.
## 1 Introduction
The existence of rapid changes (time scales from minutes to hours) in the optical brightness of blazars is a well established fact (Racine 1970; Miller et al. 1989; Romero et al. 1999; Stalin et al. 2004, 2005). These variations are usually called *microvariability* or *intranight* variability. The incidence of the phenomenon on different classes of active galactic nuclei (AGNs) seems to be also different (Jang & Miller 1997; Romero et al. 1999, 2002; Stalin et al. 2005). BL Lacertae objects, along with flat radio-spectrum quasars, seem to be among the most variable sources on short time scales. BL Lacs, in turn, can be divided into two groups: X-ray selected BL Lacs (XBLs) and radio-selected BL Lacs (RBLs) according to their spectral energy distributions (SEDs). In general, both types of objects have SEDs with two peaks, one thought to be due to synchrotron emission and the other produced by inverse Compton upscattering of lower energy photons. In the case of XBLs, the synchrotron peak falls in the X-ray band. In RBLs this peak is shifted towards lower energies, being between the radio and infrared bands. The optical microvariability behaviour of both sub-types of BL Lacs seems to be quite different (Heidt & Wagner 1996, 1998; Romero et al. 2002). The XBLs are usually less variable, with smaller duty cycles and variability amplitudes than the RBLs.
Although the optical microvariability phenomenon has been extensively studied for different AGNs, little is known about the polarimetric behaviour of these objects on time scales of hours. High-temporal resolution polarimetry has been performed for a few particular sources (e.g. 3C279, Andruchow et al. 2003) confirming the existence of microvariability in the optical polarization of some objects, but there is still no statistically significant information.
In this paper we present, by first time, results of a systematic search for microvariability in the polarization of a sample of BL Lac objects. We have looked for rapid changes in the degree of linear polarization and in the position angle of both XBLs and RBLs. Our total sample has 18 objects, so we will be able to extract some statistical conclusions on the duty cycles for these sub-types of sources. For most of the objects, this is the first time that high time resolution polarimetric observations are performed.
The structure of the paper is as follows. In Section 2 we present a detail of the observed sample. In Section 3 we describe the polarimetric observations. In Section 4.1 we describe the statistical data analysis. Section 4.2 presents individual notes on the behaviour of each source. In Section 4.3 we briefly comment on the normalized Stokes parameters. In Section 5 we discuss the statistics and the origin of the observed variability. Finally, we close with our conclusions in Section 6.
## 2 Sample
The sample adopted for this work consists of 18 sources: 10 radio-selected blazars (RBL) and 8 X-ray-selected blazars (XBL), and is taken from the AGN catalogs by Vรฉron-Cetty & Vรฉron (1998) and Padovani & Giommi (1995). We selected blazars with declinations lower than $`+20^{}`$ and brighter (at the time of our observations) than magnitude $`V=18.5`$. The redshifts spanned the range from $`z=0.044`$ to $`z=1.048`$. Basic general information on these objects is given in Table 1. In this table, Column (1) gives the name of the sources, Columns (2) and (3) give their equatorial coordinates, Column (4) their classification, Columns (5), (6) and (7) provide the colour excesses, the redshifts and the visual magnitudes taken from the NASA Extragalactic Database (NED), Column (8) gives the published maximum degree of optical linear polarization (for those cases in which this value was know from the literature), and Column (9) gives the corresponding references for Column 8.
## 3 Polarimetric observations and data reduction
The observations were done with the 2.15-m Jorge Sahade telescope at CASLEO, San Juan, Argentina, during 22 nights in April and November 2002, May 2003, and April 2004. In all occasions, we used the CASPROF photopolarimeter. This is an instrument developed at CASLEO and based on other, similar two-channel photopolarimeters, as MINIPOL and VATPOL (Magalhรฃes et al. 1984; Martรญnez et al. 1990). The observations were carried out using always the same configuration: a Johnson $`V`$ filter and an 11.3 arcsec aperture diaphragm. Integration times varied between 300 and 900 s, depending on the object brightness and the quality of the night. In all cases, we observed the target followed by a sky integration. Standard stars chosen from the catalog by Turnshek et al. (1990) were observed to determine the zero point for the position angle and the instrumental polarization; the latter was found to be practically zero. Weather conditions varied along the whole campaign, from photometric to partially cloudy (thin cirrus).
The data were processed using a systematic method, after discarding some data points affected by moonlight contamination or passing clouds. We averaged each two consecutive target observations (on the $`QU`$ plane) in order to improve the signal-to-noise ratio. A factor that affects both the object and the sky, is the presence of the Moon above the horizon. However, any systematic error, leading to spurious variations in the sky polarization, should be removed when the data are reduced, because each sky observation is made near in time and position to the corresponding source measurement. To prevent against errors due to rapid sky variations, we interpolated the sky flux and polarization (on the $`QU`$ plane) to the time corresponding to each object observation, thus giving a more accurate sky subtraction.
For each observing session we searched for any residual systematic errors by plotting the sky magnitude, as well as its polarization percentage and position angle, against time, and then comparing these graphs with the corresponding time-curves for the sources. No spurious variations, due to rapid sky changes, were evident. Two examples of typical microvariability curves are shown in Figures 1 and 2; Fig. 1 presents the behaviour of the RBL object PKS 0422$`+`$004 during two consecutive nights in November 2002, whereas Fig. 2 shows the behaviour of the XBL source PKS 2155$``$304 for three nights in November 2002. Sky values are given as open symbols in both figures; note that a different scale was used for the sky plots. Both objects were always $`0.53`$ mag brighter than the sky, thus confirming that any effect due to sky variations should not be important.
We also checked for the incidence of the foreground polarization (this is the polarization generated by interstellar dust particles oriented by the magnetic field of the Galaxy). Following Hough (1996), we used the known relation $`P_{\mathrm{max}}(\%)9E_{(BV)}`$ to set an upper limit to the foreground polarization in each of our fields. In column 5 of Table 1 we give the $`E_{(BV)}`$ values taken from the NED; it results that, since we observed at relatively high Galactic latitudes, the values of $`P_{\mathrm{max}}`$ are between $`0.13\%`$ and $`1.62\%`$. With the same purpose, we observed faint stars in most of the target fields in order to check their polarization values. In all the cases, values were $`1\%`$. Thus, we are able to confirm that the foreground polarization does not affect ours results.
We have used the Stokes parameters to check whether the pattern of the observations was random in the $`QU`$ plane or not. We consider, as usual, normalized dimensionless parameters $`U/I`$ and $`Q/I`$. $`QU`$ plots for all objects of our sample are available at http://www.iar.unlp.edu.ar/garra/garra-sdata.html.
## 4 Data analysis and results
### 4.1 Statistical analysis
We quantitatively analyzed our data by computing a formal variability indicator, following the criterion of Kesteven et al. (1976), which was used by several other authors in variability studies (Altschuler 1982; Romero et al. 1994; Andruchow et al. 2003). According to this criterion, a source is classified as variable in an observing session for the observable $`S`$ if the probability of exceeding the value
$$X^2=\underset{i=1}{\overset{n}{}}ฯต_i^2(S_iS)^2$$
(1)
by chance is $`<0.1`$ %, and it is classified as non-variable if the probability is $`>0.5`$ %. In this equation, $`ฯต_i`$ is the error corresponding to each measured value $`S_i`$, and $`S`$ is the mean value of $`S`$, given by:
$$S=\frac{_{i=1}^nฯต_i^2S_i}{_{i=1}^nฯต_i^2}.$$
(2)
If the errors are random, $`X^2`$ should be distributed as $`\chi ^2`$ with $`n1`$ degrees of freedom, where $`n`$ is the number of points in the distribution.
The other parameters that quantify the variability, in amplitude as well as in timescale, are: the fluctuations index $`\mu `$,
$$\mu =100\frac{\sigma _s}{S}\%,$$
(3)
where $`\sigma _s`$ is the standard deviation of one data set; the fractional variability index of the source $`FV`$,
$$FV=\frac{S_{\mathrm{max}}S_{\mathrm{min}}}{S_{\mathrm{max}}+S_{\mathrm{min}}},$$
(4)
where $`S_{\mathrm{max}}`$ and $`S_{\mathrm{min}}`$ are the maximum and minimum values, respectively, for the polarization or the position angle. Finally, the time interval $`\mathrm{\Delta }t`$ between the extrema in the polarization curve is defined as:
$$\mathrm{\Delta }t=|t_{\mathrm{max}}t_{\mathrm{min}}|,$$
(5)
where $`t_{\mathrm{max}}`$ and $`t_{\mathrm{min}}`$ are the times when the extreme points occur.
In Tables 2 and 3 we show the values of the variability parameters for the linear polarization percentage and the position angle for the RBLs and XBLs, respectively. Column 1 gives the name of the object, Col. 2 lists the observation dates, Col. 3 shows the number of points for each night, Col. 4 gives the total duration of the observation, Col. 5 gives the mean polarization for the observing night using Equation 2, Col. 6 shows the rms $`\sigma _P`$, Col. 7 gives the value for $`FV`$, Col. 8 is $`\mathrm{\Delta }t`$, Col. 9 shows the value of $`\chi ^2`$, and Col. 10 shows the variability class (V: if the source is variable, NV: if it is not variable, and *dubious* in the cases where no definite decision could be reached using the above given criteria). Cols. 11 to 16 give the same information for the position angle.
### 4.2 Notes on individual sources
We comment now briefly on the observed behaviour of each object in our sample.
RBLs:
#### 0118$``$272:
The linear optical polarization of this object was measured by Impey & Tapia (1988) with a significantly high value ($`P_V17\%`$). This was one of the reasons for its classification as a blazar; a similar result was obtained by Mead et al. (1990). We observed this blazar on two consecutive nights, presenting variability on the first one, with the degree of polarization rising from $`P=5.77\%`$ to $`P=9.41\%`$ in about one hour. On the second night, the degree of polarization appears as not variable but with a value of about $`8\%`$. Meanwhile, the position angle was variable with similar averages on both nights.
#### 0422+004:
This blazar was reported to have quite high values of linear optical polarization, between $`722\%`$ (Angel & Stockman 1980), presenting high variability over long periods of time (months). Mead et al. (1990) observed the source on two consecutive nights, detecting a decreasing degree of polarization from the first night to the second one ($`21.4\%`$$`12.4\%`$, respectively). In Fig. 1 we show the temporal evolution of the linear polarization degree and position angle for this object, as an example of the RBL class. This is the best sampled RBL in our campaign. It resulted to be variable, both in $`P`$ and $`\theta `$, on the two nights. The degree of polarization was quite high, mostly during the second night, reaching values up to $`13\%`$. The general trend was a smooth variation within each night, with a higher mean value for the second night. These values are in agreement with the published ones. The position angle was variable; however, its mean value did not change significantly between both nights.
#### 0521$``$365:
This southern blazar has been reported to have rapid variability at radio frequencies (see Romero et al. 1995b), as well as optical flux microvariability (Romero et al. 2002). Its host was detected and classified as a luminous giant elliptical by Falomo (1994). During the present campaign, the object experienced polarization variability on the first night, but the variability is not significant on the second one. We think that the cause for this could be the relatively large error bars for the degree of polarization during that night, which could mask any variation. The typical error was about $`6\%`$ of the measurement because of the weather conditions. The mean value of the degree of linear polarization was almost the same in the two nights, about $`3\%`$. The position angle was clearly variable during the two nights, with a small but clear rotation on the second night: during the first hours, of $`5^{}/`$h in a clockwise direction, and during the last hours, of $`7.6^{}/`$h in an anticlockwise direction. Large rotations of the polarization angle have been previously found at radio wavelengths by Luna et al. (1993) for this source.
#### 0537$``$441:
This is another well-studied BL Lac object, which has shown very high optical polarization during observations carried out by Impey and Tapia in the 1980s (Impey & Tapia 1988). This object was extensively monitored by Romero et al. (2000, 2002), presenting both behaviours, as *variable* and as *non-variable* in its optical flux at different epochs. We report here that during two consecutive nights in November 2002, 0537$``$441 presented a high degree of polarization. On the first night the objectโs variability appeared as *dubious*, meanwhile on the second night it was clearly $`variable`$, undergoing rapid fluctuations with typical time scales of $`1`$ h and amplitudes of $`1.4\%`$. On the contrary, $`\theta `$ was $`variable`$ on the first night and *dubious* on the second one, with a mean $`\theta 0^{}`$.
#### 0829+046:
Its first polarization observations at optical wavelengths were presented by Wills et al. (1980), who reported variability over a few days time interval. Unfortunately, we could follow this object just one night. During this period, the object presented a very high degree of polarization (up to $`17\%`$), increasing with time, and it was variable in both the degree of polarization and position angle. It is also interesting to mention that this object was observed by Giroletti et al. (2004) at different radio wavelengths in order to resolve the jet structure and they found that it is one of the BL Lacs that presents evidence of emission on both sides of the core: two symmetric jets were detected emerging from the core and both are bended.
#### 1144$``$379:
This object was classified as a blazar by Impey & Tapia (1988); they reported values for the degree of polarization between $`0.0\%`$ and $`9.4\%`$. We observed the source in April 2002 and it was variable in both $`P`$ and $`\theta `$. During the last night, the degree of polarization showed a peculiar behaviour: it rose about $`14\%`$ in 4 hours, starting at $`P=2.5\%`$ and ending at $`P=16.5\%`$. After ruling out all possible error sources (see Sect. 3) we conclude that the cause of this peculiar behaviour is intrinsic to the source.
#### 1510$``$089:
This object was confirmed as a blazar by Moore & Stockman (1981). Previous measurements of its optical polarization degree, made in $`1980`$, were all under $`7.8\%`$; however, Mead et al. (1990) reported a high value of $`P=9.1\%`$, in the $`I`$ band. We observed 1510$``$089 during three nights, one in April 2002, and two in April 2004. The position angle was always variable. The degree of polarization was variable in 2002 and on the first night in 2004, but not on the second one. Its average value was about $`3\%`$ along all three nights, but reaching values as high as $`13.8\%`$ during $`2002`$. This BL Lac presented no microvariability during $`1998`$ and $`1999`$ in its optical flux (Romero et al. 2002).
#### 1514$``$241
(AP Lib): This is one of the objects which defined the blazar class; it has presented values of optical polarization between $`27\%`$ (Angel & Stockman 1980). Mead et al. (1990) reported similar values. We followed AP Lib on different occasions, resulting always variable; however, the average degree of polarization was quite low. During the last night in April 2004 the position angle rotated in an anti-clockwise direction from $`180.3^{}`$ to $`170.9^{}`$ with a speed of $`10.5^{}/`$h.<sup>2</sup><sup>2</sup>2Joining this information with additional data obtained without filters, the general trend in the position angle was a constant rotation with sporadic direction reversals.
#### 1749+093:
This object has displayed dramatic polarization variability at optical and infrared wavelengths (Brindle et al. 1986), whereas no significant variations were detected later (Mead et al. 1990). Typical values for the degree of the optical polarization are between $`39\%`$ (Kรผhr & Schmidt 1990). Our variability data classify the source as *dubious* during the first night and variable on the second one, with a low mean value of the degree of polarization, but reaching a maximum value of $`P_V=9.8\%`$. The position angle was variable on both nights, but this variability was probably not real, because when the modulus of the polarization vector is small, the angle is ill defined, thus preventing any real variability to be detected.
#### 2005$``$489:
As far as we know, we are presenting here the first optical polarization data for this BL Lac object. The source was variable, with a relatively high degree of polarization, during the only night when we could observe it. The position angle variability was classified as *dubious*, but in fact, it remained almost constant around $`93.7^{}`$ with a sigma of $`0.2^{}`$.
XBLs:
#### 0548$``$322:
Angel & Stockman (1980) had reported low levels in the degree of polarization ($`1.52\%`$). Similar results were obtained during the campaign undertaken by Jannuzi et al. (1993), when $`P`$ did not rise above $`4\%`$, with position angle variable and showing a $`0.5`$ mag change in its optical flux. We followed this typical BL Lac for two consecutive nights in November 2002, with a good time resolution. During both nights, the degree of optical polarization appears to be variable but low, and the position angle was variable, too.
#### 0558$``$385:
We present the first polarization results for this object. The average polarization is very low, $`P0.80.9`$ %. The object formally classifies as *variable* but, since its polarization is so low, the large fluctuations detected in the position angle might be spurious.
#### 1026$``$174:
There are no previous optical polarization data for this blazar. We observed it in May 2003, when it displayed a *dubious* behaviour in its polarization degree. However, during the second night, some variation is present, unfortunately masked out by the large error bars due to the weakness of the source. The highest and lowest values of the optical polarization were $`7.6\%`$ and $`3.9\%`$, respectively. The position angle was variable, showing no clear rotation trend, with values around $`185^{}`$.
#### 1101$``$232:
This blazar has been reported to have quite low values of optical polarization (the maximum detected was $`2.7\%`$), with evidences of intrinsic variability (Jannuzi et al. 1993). We have observed this source in 2002, 2003 and 2004. The behaviour displayed by the source went through different stages (from $`V`$ to $`NV`$) along the different opportunities we had to observe it. During the only night that we observed it in May 2003, $`P_V`$ rose to $`14.7\%`$, a very high value for this object, which had previously presented lower polarization values.
#### 1312$``$422
: This is another BL Lac with no previous data on its optical polarization. We just followed it on one night in April 2002 and another night in April 2004. The degree of polarization was quite low ($`P_V3.4\%`$) and not variable on both occasions.
#### 1440+122:
Published information about this object is scarce. Recent radio observations (Giroletti et al. 2004) revealed more details about its structure, but no previous optical polarization data are available. This blazar resulted to be variable in both $`P_V`$ and $`\theta `$ during April 2002. $`P_V`$ reached as high a value as $`8.3\%`$, with a peculiar behaviour during the first night we followed it. After discarding any possible error sources (see Sect. 3), the trend in the polarization degree is an interesting one, rising from $`P_V=1.3\%`$ to $`P_V=8.3\%`$ during the first $`1.8`$ h, then going down to $`P_V=0.5\%`$ in about $`1`$ h, and finally rising again up to $`P_V=7.5\%`$. Meanwhile, the mean polarization was not too high ($`P_V=3.4\%`$). On the second night, $`P_V`$ was variable, but with no peculiar trend. The position angle did not follow the variations in $`P_V`$; however, it was always variable, with values around $`\theta 95^{}`$.
#### 1553+133:
The degree of polarization of this BL Lac object was variable during our observations, with a flickering behaviour on the first night. A qualitatively similar flickering has been detected before at radio wavelengths in extragalactic radio sources (Heeschen 1984; Romero et al. 1995a). In the optical polarization the origin of the rapid flickering must be intrinsic to the source, probably related with turbulence in the magnetic field of the inner jet. The upper limit of the polarization was $`4.2\%`$ for this object.
#### 2155-304:
This is a very well studied BL Lac object, known for its short variability time scales at optical to X-ray wavelengths (Jannuzi et al. 1993). Angel & Stockman (1980) reported values for the degree of optical polarization between $`37\%`$. In the mid 1980s, Brindle et al. (1986) detected polarization percentage variations and a clockwise rotation of the position angle, although on a $`48`$ h time scale. Four years later, Mead et al. (1990) measured a higher than usual polarization ($`P_V10\%`$). The extensive monitoring made by Smith et al. (1992) in optical polarization and photometry revealed clear variations over long time scales; the authors also reported a $`2\%3\%`$ variation in $`P_V`$ and as much as $`25^{}`$ in $`\theta `$, in $`24`$ h. They also found a mild wavelength dependent polarization and a more rapid variation in the optical polarization than in the total optical flux. In a more recent work Tommasi et al. (2001) reported a multiband monitoring of the optical polarization searching for intranight and also long term variability. This campaign was made using the 2.15 m telescope at CASLEO equipped with the Photopolarimeter of the Turin Observatory and the total observing time was $`47`$ h along four different periods in 1998 and 1999. The authors reported a $`P_V`$ lower than $`7\%`$ with small amplitude intranight variations, $`1.3\%`$ in $`P`$ and $`7^{}`$ in $`\theta `$, but no statistical analysis of this behaviour was reported. They also found a low wavelength dependence in both the linear polarization and the position angle.
Being a relatively bright object, we were able to use short exposure times, hence getting well-sampled time curves. These are shown, as an example for the XBL class, in Fig.2. Significant variations both in $`P_V`$ and $`\theta `$ are clearly seen on each of the three consecutive nights that we followed the source, with a moderately high mean polarization $`P_V5\%`$ and with a position angle varying from $`\theta =87^{}`$ to $`\theta =105^{}`$. During the first night, $`P_V`$ raised as high as $`5.7\%`$ at the beginning of the night and then went down, ending at $`4.7\%`$. Apparently, this decreasing trend continued during the day hours, because at the beginning of the second night, $`P_V`$ started at $`3.2\%`$, going up for the rest of the night and presenting an inverse behaviour during the last night. With respect to the position angle, the variation was clear and presented a fast rotation during the third night. The angle rotated in an anti-clockwise direction from $`96.7^{}`$ to $`93.5^{}`$ with a speed of $`1.9^{}/\mathrm{h}`$.
## 5 Discussion
In order to characterize the two different classes of objects under study here, we analyzed the behaviour of the sources from a statistical point of view. First, we calculated the duty cycles (DC) for the sources of a given class. This quantity can be estimated, following Romero et al. (1999, 2002), as
$$DC=100\frac{_{i=1}^nN_i,(1/\mathrm{\Delta }t_i)}{_{i=1}^n(1/\mathrm{\Delta }t_i)}\%,$$
(6)
where $`\mathrm{\Delta }t_i=\mathrm{\Delta }t_{i,\mathrm{obs}}(1+z)^1`$ is the duration, corrected by the corresponding redshift, of the $`i`$-th data set of the quantity and class under study; $`N_i`$ is the weight (equal to 1 if the source was classified as $`V`$, or 0 if the source was NV or *dubious*). Because we weighted *dubious* cases with 0, the DCs calculated are actually lower limits. The corresponding DCs for the degree of polarization ($`P`$) and position angle ($`\theta `$) for both classes of objects (RBLs and XBLs) are: $`DC(P,\mathrm{RBL})=77.01\%`$, $`DC(\theta ,\mathrm{RBL})=87.25\%`$, and $`DC(P,\mathrm{XBL})=51.23\%`$, $`DC(\theta ,\mathrm{XBL})=55.15\%`$. So, the RBLs appear to constitute the most variable class. A similar behaviour has been found when only optical flux variations were considered by Romero et al. (2002), with duty cycles $`DC=71.5\%`$ and $`DC=27.9\%`$ for the RBLs and XBLs, respectively. The photometric microvariability of XBLs seems to be systematically lower, nonetheless, than their polarization microvariability.
A complementary view can be obtained by plotting the histograms of the distributions of the sources that resulted to be variable. Figures 3 and 4 show the number of sources classified as *variable* against the time scales for the variation (Col. 8 and 14 in Tables 2 and 3), for the degree of polarization and position angle, respectively.
It can be seen that the RBLs have a wider and flatter distribution in the degree of polarization than the XBLs; this is an indicator that, when an XBLs is variable, its variation timescale is shorter than that of the RBLs (typically $`\mathrm{\Delta }t1`$ h). The histograms corresponding to the position angles present no significant differences between the two classes.
Similar histograms were made showing the distributions of variable sources against the fractional variability index (Col. 7 and 13 in the same tables). This is shown in Figs. 5 and 6 for the same parameters as before. Here, the RBLs appear to have two peaks, one for $`P0.20.3\%`$, and the other around $`P0.80.9`$. On the contrary, the XBLs have a more uniform distribution.
Concerning the position angle, the variation of the RBLs appears to be more frequent for the lowest $`FV`$ values. Again, the XBLs seem to have a more uniform distribution.
Since we present here a significant number of sources belonging to two different sub-types of BL Lac objects, it is also interesting to compare the average degree of optical polarization between them. Our results confirm the previous ones reported by Fan et al. (1997): the XBLs generally have lower optical polarization than the RBLs; this kind of behaviour can be seen in the histograms drawn in Figure 7. A Kolmogorov-Smirnov test shows that both data sets are most probably taken from different parent distributions, although the significance level is only moderately high (95%). As an additional information, we also include the distribution of V and NV or *dubious* cases for both classes of BL Lacs studied here.
In general, since XBLs have the synchrotron peak of their spectral energy distribution at X-ray energies, we could expect that these objects have, on average, either higher magnetic fields or more energetic particles than RBLs, which peak at radio-IR wavelengths. The fact that they have on average less polarization and that this polarization is less variable than what is found for RBLs seems to support the second possibility, i.e., the particles in their jets are systematically more energetic than in RBLs. On the contrary, as noticed by Fan et al. (1997), the RBLs seem to have larger macroscopic relativistic motions, hence displaying higher duty cycles for rapid variability, which is probably associated with relativistic shocks in the jets. Their magnetic fields seem to be also systematically stronger than in XBLs, as indicated by the higher degree of linear polarization. This leads to a simple picture where XBLs have particles with high microscopic Lorentz factors that cool radiating high-energy synchrotron emission whereas RBLs have less energetic particles but higher macroscopic bulk motions and stronger fields, hence presenting higher variability. Alternatively, XBL could have similar magnetic fields, but less homogeneous, hence less degree of linear polarization. The origin of the rapid microvariability seems to be associated with relativistic shocks in any case (e.g. Romero et al. 1995b, and references therein).
## 6 Conclusions
We have monitored 8 XBLs and 10 RBLs, looking for intranight variability in the optical polarization. We have found high duty cycles for both the degree of linear polarization and the polarization angle of RBLs. The average polarization is also stronger than for XBLs. XBLs, although displaying a lower level of polarimetric microvariability, are also significantly variable, with duty cycles of $`50`$ %, higher than what is observed from purely photometric observations. We speculate, on the basis of our findings, that the stronger synchrotron losses presented by XBLs might be due to systematically higher microscopic Lorentz factors for the particles in the jets, rather than to stronger magnetic fields. However, it should be noted that further observations of objects (with each object monitored on several nights) are needed to establish this conclusion firmly.
###### Acknowledgements.
This work has been supported by CONICET and ANPCyT (through Grant PICT 03-13291). We thank the staff of CASLEO observatory for valuable help during the observations. We also thank the constructive suggestions made by the anonymous referee. This research made use of the NASA Extragalactic Data system.
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# Statics and Dynamics of Vortex Liquid Crystals
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## Abstract
Using numerical simulations we examine the static and dynamic properties of the recently proposed vortex liquid crystal state. We confirm the existence of a smectic-A phase in the absence of pinning. Quenched disorder can induce a smectic state even at $`T=0`$. When an external drive is applied, a variety of anisotropic dynamical flow states with distinct voltage signatures occur, including elastic depinning in the hard direction and plastic depinning in the easy direction. We disuses the implications of the anisotropic transport for other systems which exhibit depinning phenomena, such as stripes and electron liquid crystals.
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Recently, a new state of vortex matter termed a vortex liquid crystal was proposed to occur in superconductors with anisotropic vortex-vortex interactions . In such systems, the vortex lattice first melts in the soft direction, giving rise to an intermediate vortex smectic-A state, followed at higher temperatures by a melting into a nematic state. The initial theoretical calculations of this transition combined an elastic model with the Lindemann criterion for melting; however, the validity of this approach has been called into question . The Lindemann criterion does not take into account the proliferation of dislocations that is likely to occur in a smectic-A state, so a numerical investigation would be very useful both to determine whether the smectic-A state can occur in the vortex system as well as to examine the dynamics of vortex liquid crystals in the presence of disorder. The physics of the vortex liquid crystal state should be generic to the class of problems which can effectively be modeled as a two-dimensional (2D) system of particles with anisotropic repulsive interactions. Such a system has already been physically realized for magnetic colloidal particles, where a smectic-A state was observed along with dislocations that have preferentially aligned Burgers vectors . There is also considerable interest in electron liquid crystal states, which may arise when anisotropic interactions in classical electron crystals give rise to smectic and nematic states . Evidence for such states has been observed in transport measurements which show hard and soft directions for flow .
Previous studies of vortex smectic states considered vortices interacting with some form of an underlying 1D periodically modulated substrate . A very similar system in which smectic states have been observed is colloidal particles interacting with 1D periodic substrates . Each of these systems melts into an intermediate smectic-C state. When the underlying substrate is disordered rather than periodic, application of an external drive induces an anisotropic fluctuating force that organizes the vortices into a moving smectic state, where the dislocations in the vortex lattice are aligned with the direction of the applied drive . In the case of the proposed vortex liquid crystal state, the anisotropy arises when the vortex cross section becomes elliptical due to an anisotropic superfluid stiffness which leads to different effective masses in the three crystalline directions . The theoretical calculations in Ref. were performed for a system with no quenched disorder; however, real superconductors often contain significant amounts of random pinning. It would be desirable to understand the transport properties of vortex liquid crystals in order to identify signatures of the liquid crystal phase and seek new types of dynamical phenomena in these systems. Understanding how quenched disorder affects an anisotropic system of repulsively interacting particles is also relevant to dynamics in electron liquid crystal states.
To address these issues, we consider a 2D system of $`N_v`$ interacting vortices with periodic boundary conditions in the $`x`$ and $`y`$ directions. The overdamped equation of motion for a single vortex $`i`$ is
$$\eta \frac{d๐_i}{dt}=๐_i^{vv}+๐_i^T+๐_i^p+๐_i^d$$
(1)
The damping constant $`\eta `$ is set to unity. The vortex-vortex interaction force is $`๐_i^{vv}=_{ji}^{N_v}A_vK_1(r_{ij}/\lambda )\widehat{๐ซ}_{\mathrm{๐ข๐ฃ}}`$, where $`K_1`$ is the modified Bessel function, which decays exponentially for large distances, $`\lambda `$ is the London penetration depth, $`A_v`$ is the vortex interaction prefactor, and $`r_{ij}`$ is the distance between vortices $`i`$ and $`j`$. The Bessel function is appropriate for stiff, 3D vortex lines. We have also considered $`1/r`$ interaction potentials appropriate for thin film superconductors as well as Yukawa interaction potentials for colloidal particles and find the same qualitative features. The thermal force $`๐_i^T`$ arises from random Langevin kicks with the properties $`<๐_i^T>=0`$ and $`<๐_i^T(t)๐_j^T(t^{})>=2\eta k_BT\delta (tt^{})\delta _{ij}`$. The quenched disorder $`๐_i^p`$ is modeled as random pinning sites in the form of attractive parabolic traps of radius $`r_p=0.2\lambda `$ and strength $`f_p`$. The Lorentz driving force from an external applied current is $`๐^d`$. The system size is measured in units of $`\lambda `$ and the forces in terms of $`A_v`$. The anisotropic interactions are introduced by multiplying the vortex-vortex interaction force in the $`x`$ and $`y`$ directions by a vector $`(C_x,C_y)`$, where the anisotropy $`C=C_x/C_y`$. In this work we concentrate on the case $`C=1/\sqrt{10}`$ which
is the value considered in Ref. . We take the $`x`$ axis to be the soft direction and the $`y`$ axis as the hard direction.
We first consider the case where the pinning and the external driving force are absent. In Fig. 1 we illustrate the melting of a $`24\lambda \times 24\lambda `$ system with a vortex density of $`\rho _v=1.2/\lambda ^2`$. Figure 1(a) shows the vortex positions (dots) and trajectories (lines) for a fixed period of time with a fixed $`T=0.5`$, and Fig. 1(b) shows a corresponding Delaunay triangulation. At this temperature, the system remains in a crystalline state with no dislocations. The vortices are undergoing larger random displacements in the soft ($`x`$) direction than in the hard ($`y`$) direction; however, there is no long time diffusion of the particles. Figures 1(c) and 1(d) present the smectic-A state at $`T=1.2`$. Here the trajectories have a 1D liquid structure with
motion along the soft $`x`$ direction and no significant translation of the vortices in the $`y`$ direction. The Delaunay triangulation indicates the presence of dislocations with aligned Burgers vectors, which is characteristic of the smectic-A state. Figures 1(e) and 1(f) illustrate the vortex liquid phase at $`T=1.35`$. The vortex trajectories show clear diffusion in both the $`x`$ and $`y`$ directions, with more pronounced motion in the $`x`$ direction. The dislocations are no longer aligned in a single direction, indicating the loss of long-range order in both the $`x`$ and $`y`$ directions. These results confirm that a smectic-A state can occur in a system of vortices with anisotropic interactions, as predicted by theoretical calculations . We note that when the anisotropy ratio $`C`$ is too small, the two-step melting transition illustrated here is lost.
To further characterize the smectic state, in Fig. 2 we plot the average particle displacements for the $`x`$ and $`y`$ directions, $`d_x=<_i^{N_v}|x_i(t)x_i(t^{})|>/N_v`$ and $`d_y=<_i^{N_v}|y_i(t)y_i(t^{})|>/N_v`$. In the smectic phase at $`T=1.21`$, shown in Fig. 2(a), $`d_x/a`$ increases much more rapidly than $`d_y/a`$, and does not saturate but increases to a value over 1, indicating that the vortices can diffuse more than a lattice constant in the $`x`$ direction over time. This is due to the formation of dislocations which allow adjacent rows of vortices to slip past each other while remaining confined in the $`y`$ direction. We note that the saturation value of $`d_y/a`$ is approximately $`1/5`$, larger than the Lindemann criterion value of $`1/10`$. Excess motion in the $`y`$ direction occurs during a sliding event when two rows slip past each other and the vortices in each row are temporarily displaced in the direction perpendicular to the slip plane. This transverse motion is not large enough to permit the formation of dislocations aligned in the hard direction. In the nematic phase, shown in Fig. 2(b) at $`T=1.35`$, $`d_x`$ still increases more rapidly than $`d_y`$; however, the continuous increase of both quantities indicates that the particles are diffusing throughout the entire system.
We next consider the effect of random disorder by adding $`N_p=2N_v`$ randomly located pinning sites to the same system studied in Fig. 1, and then conducting a
series of simulations at varied $`T`$ and varied pinning strength $`f_p`$. For high temperatures we always obtain a nematic phase, while at $`T=0`$ for large $`f_p`$ we find a pinned nematic phase. At low $`T`$ and low $`f_p`$ we observe a phase very similar to that shown in Fig. 1(c,d), where the vortex lattice is oriented in the soft direction and there are a small number of aligned dislocations. We term this a pinned smectic-A phase. In Fig. 3 we indicate the regions in which the smectic and disordered phases appear as a function of temperature and pinning strength. The phase boundary is identified via the density of sixfold coordinated particles, $`P_6`$; the defect density is given by 1-$`P_6`$. In the crystal phase, there are no defects and $`P_6=1`$. In the smectic phase, $`P_6=0.91`$ to $`0.95`$, and in the nematic phase $`P_6>0.8`$. In the inset of Fig. 3 we plot $`P_6`$ vs $`T`$ for two different disorder strengths. For $`f_p=0.025`$ (upper line) the system is in the pinned smectic state at $`T=0`$. As $`T`$ increases, there is a clear transition to the disordered state, as indicated by the drop in in $`P_6`$ near $`T=1.19`$. The lower line shows $`P_6`$ for $`f_p=0.2`$, when the pinning is strong enough to disorder the system even at $`T=0`$. These results suggest that weak random disorder can increase the extent of the regions where the smectic-A phase occurs when there are anisotropic interactions, by suppressing the crystalline phase at low temperatures and raising the melting temperature of the smectic state.
It has been shown that dislocations are always present for weak disorder in two dimensional isotropic systems; however, the distance between the dislocations can be arbitrarily large compared to the range of the
translational order, so that for a wide range of temperatures and disorder strengths the system behaves as a 2D Bragg glass . In the vortex liquid crystal case there are two length scales associated with the hard and soft directions, and the disorder induced dislocations first form in the soft direction. It may be possible that, on very large length scales, dislocations in the hard direction will also appear. For the parameters considered here, dislocations are present except at the lowest pinning strengths, roughly indicated by a dashed line in Fig. 3, where a 2D anisotropic Bragg glass forms.
We next consider dynamical effects in the presence of pinning. In the smectic-A state, there should be distinct transport signatures for the hard and soft directions. When the disorder is strong enough to destroy the smectic phase, there may still be an anisotropic transport signature if the system retains some form of nematic order. We first consider the pinned smectic-A state found at $`f_p=0.04`$ and $`T=0`$. We perform separate simulations for driving in the soft direction, $`๐_d=f_d\widehat{๐ฑ}`$, and the hard direction, $`๐_d=f_d\widehat{๐ฒ}`$, increasing the applied drive very slowly to avoid any transient effects. In Fig. 4(b) we plot $`V_y=(1/N_v)<_i^{N_v}v_y>`$ (upper curve) for driving in the hard direction and $`V_x=(1/N_v)<_i^{N_v}v_x>`$ (lower curve) for driving in the soft direction. Here, $`V_x<V_y`$, indicating that motion in the soft direction is easier except at very high drives when the effects of the pinning are washed out and the two curves come together. In Fig. 4(a) we plot $`P_6`$ for the two different driving directions. At $`f_d=0,`$ $`P_6`$ is slightly less than one due to the presence of a small number of dislocations in the smectic state. For depinning in the hard direction, $`P_6^y`$ (upper curve), the system depins elastically without a proliferation of defects, and the vortices do not exchange neighbors as they move. For driving in the easy direction, $`P_6^x`$ (lower curve) drops substantially when the vortices depin plastically, and a portion of the vortices remain pinned while others flow past. The effective pinning is known to be higher for a soft system where defects can proliferate than for an elastic system . A similar proliferation of defects has been associated with the so called peak effect, where there is a sudden increase in the effective pinning force as a function of temperature or applied magnetic field . At high drive, the system shifts from plastic flow to a dynamically reordered state as indicated by the increase in $`P_6^x`$ in Fig. 4(a), as well as by the merging of $`V_x`$ and $`V_y`$ in Fig. 4(b). These results imply that, in the pinned smectic-A state, the depinning is elastic in the hard direction and plastic in the soft direction. Further, the elastic and plastic depinning transitions produce different scaling responses in the velocity force curves. At depinning, the velocity scales with the driving force in the form $`V=(f_df_c)^\beta `$ . For the plastic flow regime we find $`\beta >1.0`$ while in the elastic flow regime we find $`\beta <1.0`$, in agreement with theoretical expectations.
At finite temperatures and for $`f_p`$ large enough that we observe only plastic depinning in both directions, we observe that the critical depinning force in the soft direction, $`f_c^x`$, is lower than the depinning force in the hard direction, $`f_c^y`$, even though $`V_x<V_y`$ at intermediate drives. This implies that the anisotropic flow exhibits a reversal from $`V_x>V_y`$ to $`V_x<V_y`$ at low drives. We explicitly demonstrate this effect for a system with $`f_p=0.25`$ and $`T=0.25`$ in Fig. 4(c,d). Here, the depinning is plastic in both directions, and the depinning forces are $`f_c^x=0.015`$ and $`f_c^y=0.11`$. There are fewer dislocations for $`๐_d=f_d\widehat{๐ฒ}`$ and the system reorders at $`f_d^y=0.2`$. For $`๐_d=f_d\widehat{๐ฑ}`$, the system does not reorder until $`f_d^x=0.8`$. There is a clear crossing of the velocity force curves at $`f_d=0.13`$ so that the flow is easier in the soft direction for $`f_d<0.13`$ and easier in the hard direction for $`f_d>0.13`$. In Fig. 4(d) we show a blowup of this region. The crossing of the velocity force curves can be understood by considering that the depinning in the soft direction is plastic. At low drives, individual vortices can be thermally activated, giving rise to creep. For driving in the hard direction, the depinning is elastic and individual vortex hopping is not possible, so that only collective creep can occur. In the case of the nematic phase, when there is some plastic flow in the $`y`$-direction, there is still a large correlated length scale that must move so thermal effects are greatly reduced. Thus, creep in the pinned smectic phase and pinned nematic phase is enhanced in the soft direction compared to the hard direction. Even for very high values of $`f_p`$ we observe an intermediate anisotropic response, suggesting the system can be considered a pinned nematic phase.
In conclusion, we have performed simulations of the recently proposed vortex liquid crystal state where the vortex-vortex interactions are anisotropic. We find that, in the absence of disorder, the system shows an intermediate melting into a smectic-A state as proposed theoretically in Ref. . The smectic-A state contains a small fraction of dislocations which are all aligned in the soft direction. In the presence of disorder, a pinned smectic-A state can occur, and the system depins plastically in the soft direction but elastically in the hard direction. We predict that, for equal intermediate drives, the velocity in the soft direction will be lower than in the hard direction. At finite temperatures the creep is much higher in the soft direction due to the fact that individual vortex hopping can occur, whereas creep is suppressed in the hard direction since the vortex motion is much more correlated. For high temperatures and disorder strengths, the system is disordered. For strong disorder the anisotropic transport should still be observable in the pinned nematic phase. We note that many of our results may apply to electron liquid crystals as well.
We thank E. Carlson for useful discussions. This work was supported by the U.S. Department of Energy under Contract No. W-7405-ENG-36.
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# First quantized approaches to neutrino oscillations and second quantization
## I Introduction
Compelling experimental evidences exp have shown that neutrinos undergo flavor oscillations in vacuum. Consequently, this fact requires massive neutrinos with mixing. These ingredients are not present in the standard model of elementary particles. For this reason, on the one hand, neutrino oscillations can provide a direct window to probe physics beyond the standard model Smirnov:03 . On the other hand, some theoretical studies of mixing in the context of quantum field theory (QFT) by Blasone and Vitiello (BV) BV:AP95 ; BlasoneP:03 show the mixing problem may be related to more fundamental issues such as unitarily inequivalent representations and the vacuum structure, and its study is theoretically interesting for its own sake.
Nevertheless, the simpler plane-wave quantum mechanical descriptions Pontecorvo ; PDG seemed to be in accordance, in certain realistic limits, with more refined descriptions, including various ingredients, such as localization aspects Kayser:81 ; Giunti:91 ; ccnishi04 , flavor current densities Zralek:98 , influence of creation and detection processes KiersWeiss ; Giunti:wp:coh , time-dependent perturbation theory Rich , and intermediate neutrinos with path integrals Field . Moreover, many treatmens within the quantum field theory (QFT) framework were also proposed Giunti:93 ; Campagne ; Beuthe:PRep ; Kobzarev ; GS ; Dolgov ; Cardall.00 , aiming to solve the various unclear aspects of the quantum mechanics of neutrino oscillations Rich ; Zralek:98 .
It has been known for a long time that the coherence necessary for neutrino oscillations depends crucially on localization aspects of the particles involved in the production of neutrinos Kayser:81 . This point of view can be supported by QFT arguments GS as well. It raises then the question of how the coherent superposition of mass eigenstate neutrinos, which is called a โflavorโ eigenstate, is created Dolgov . One way that became customary to avoid the ambiguities involving the question on how neutrinos are created and detected is to use an external (E) wave packet (WP) approach Beuthe:PRep , in contrast to an intermediate (I) WP approach. According to Ref. Beuthe:PRep , the IWP treatments are the simpler first quantized ones treating the propagation of neutrinos as free localized wave packets. In contrast, EWP approaches consider localized wave packets for the sources and detection particles while the neutrinos were considered intermediate virtual particles. The central issue distinguishing the general IWP and EWP approaches is: despite its direct unobservability, is the intermediate neutrino a real (on-shell) particle propagating freely? If the answer is affirmative, the IWP approaches would be a good approximation of the oscillation phenomena.
On the other hand, another classification scheme can be used do classify the various existing treatments considering a more physical criterion irrespective of the use of WPs. Consider the descriptions of neutrino oscillations that (A) include explicitly the interactions responsible for the mixing and those (B) that only treat the propagation of neutrinos, i.e,. the mixing is an ad hoc ingredient. A more subtle aspect in between would be the (explicit or phenomenologically modelled) consideration of the production (and detection) process(es). In general, the IWP approaches are of type (B). The EWP approaches are of type (A). The BV approach, although in the QFT formalism, is of type (B) since mixing is introduced without explicitly including the interaction responsible for it. The type (B) approaches have the virtue that they can be formulated in a way in which total oscillation probability in time is always conserved and normalized to one ccnishi04 ; BV:AP95 . This feature will be present in all first quantized approaches treated here (secs. II and III) and in a second quantized version (sec. IV.1). If different observables are considered, or a modeling of the details of the production and detection processes is attempted, further normalization is necessary Giunti:91 ; BlasoneP:03 ; Giunti:wp:coh . In such cases, the oscillating observable might differ from the oscillation probability. On the other hand, type (A) approaches tend to be more realistic and can account for the production and detection processes giving experimentally observable oscillation probabilities Cardall.00 . Of course, they are essential to the investigation of how neutrinos are produced and detected KiersWeiss ; Dolgov . We are not directly interested in these matters here.
Considering first quantized type (B) approaches, some recent works treating the flavor oscillation for spin one-half particles Bernardini:Euro ; Bernardini:PRD05 have already find additional oscillatorial effects compared to usual oscillation formulas with WPs Giunti:91 ; ccnishi04 . These effects are investigated and it is shown in sec. II how these additional oscillatorial behavior, which have characteristic frequencies much greater than usual oscillation frequencies, comes from the interference between positive and negative frequency components of the initial WP. It can be understood as a consequence of the impossibility to simultaneously exclude all negative energy contributions of the initial spinorial wave function with respect to bases characterized by different masses. Moreover, this rapid oscillations are always present, independently of the initial WP, if a well defined flavor is attributed to the initial WP.
To make clear the origin of the additional oscillatory contributions, we calculate, in sec. III, the oscillation formula for a charged spin 0 particle in the Sakata-Taketani Hamiltonian formalism FV , which is equivalent to the Klein-Gordon scalar wave equation. (The explicit analysis with mixed Klein-Gordon equation is made in Ref. Dvornikov, , paying special attention to the relativistic initial value problem.) The oscillation formula in this case also possesses the additional interference terms between positive and negative frequency parts, very similar to the one obtained in the spin 1/2 case. From this example we will see that these interference terms are inevitable from a relativistic classical field theory perspective where covariance and causality is required. It is not specially associated to the spin degree of freedom.
Another particular ingredient of neutrino oscillations can be included naturally within Dirac theory: the left-handedness of neutrinos created and detected through weak interactions. This fact, for a Dirac neutrino Mohapatra , implies an additional probability loss due to conversion of left-handed neutrinos into right-handed neutrinos, which is possible because chirality is no longer a constant of motion for massive Dirac particles chiral . Although previous calculations Bernardini:PRD05 have shown an approximate contribution to this effect, we calculate in sec. II.1 the complete effect.
Concerning type (A) approaches, specifically the EWP description, we are interested to analyze further how is the propagation of intermediate virtual neutrinos. The framework where the investigations on first quantized approaches are made here is based on the calculation of the evolution kernels for free theories in presence of mixing. This enable us to deduce general oscillation probabilities in which there is explicit decoupling from the oscillating part (where all the oscillation information rests) and the initial wave packet. Another advantage of doing the calculations this way is that it resembles the propagator methods in covariant perturbation theory, which EWP approaches are based on. The free evolution kernel for fermions have a close relationship with the Feynman propagator used in QFT. What is common to both is that both particle and antiparticle parts contribute to the evolution or propagation. The necessity of the negative frequency part in the free evolution kernel is required from completeness and causality arguments but it also leads to the interference of positive and negative frequencies in flavor oscillation, treated in secs. II and III. Then the question also arises in EWP approaches: are there contributions from both particles and antiparticles in the propagation of virtual neutrinos? In a simple microscopic scattering process, this question is meaningless since virtual particles are usually off-shell particles and must naturally have both contributions. However in EWP approaches the neutrinos propagate through macroscopic distances and, indeed, it can be shown GS ; Dolgov that the virtual neutrinos are on-shell particles. There is no discussion, though, about the possibility of neutrino and antineutrino contributions to the process; both can be on-shell. This investigation is carried on in sec. IV calculating explicitly the amplitude of production/propagation/detection process in an EWP approach.
As a last task, we develop a simple, type (B), second quantized description of flavor oscillation in sec. IV.1 using the free second quantized spin 1/2 fermionic theory in presence of mixing. This treatment has some similarities with the BV formalism but it does not require the introduction of flavor Fock spaces and Bogoliubov transformations. It means that the Fock space considered will be the one spanned by the mass eigenstates. Within this formalism it will be shown that the additional rapid oscillation contributions calculated through first quantized approaches do not survive the second quantization since only superpositions of particles (antiparticles) are used as initial neutrino (antineutrino) โflavorโ states. Moreover, this property is not satisfied in the BV approach because the BV flavor states are mixtures of particle and antiparticle components; this is the ingredient responsible for a different oscillation probability BV:AP95 .
## II Flavor Oscillation for Dirac fermions
It is well known that the Dirac equation can give a significantly good description of a Dirac fermion if its inherent localization is much bigger than its Compton wave length; usually this is associated with weak external fields. For example, the spectrum for the hydrogen atom can be obtained with the relativistic corrections included (fine structure) (IZ, , p. 72). One of the terms responsible for fine structure, the Darwin term, can be interpreted as coming from the interference between positive and negative frequency parts (zitterbewegung) of the hydrogen eigenfunction in Dirac theory compared to the nonrelativistic theory FV . On the other hand a situation where the theory fails to give a satisfactory physical description is exemplified by the Klein paradox (IZ, , p. 62): the transmission coefficient for a electron moving towards a step barrier becomes negative for certain barrier heights, exactly when the localization of the electron wave function inside the barrier is comparable with its Compton wave length.
Bearing in mind that first quantized approaches may fail under certain conditions we will treat in this section the flavor oscillation problem using the free Dirac theory in presence of two families mixing. The extension to treat three families of neutrinos is straightforward. A matricial notation will be used throughout the article for the first quantized approaches to express the mixing.
In matricial notation the mixing relation between flavor wave functions $`\mathrm{\Psi }_f^๐ณ(๐ฑ)(\psi _{\nu _e}^๐ณ(๐ฑ),\psi _{\nu _\mu }^๐ณ(๐ฑ))`$ and mass wave functions $`\mathrm{\Psi }_m^๐ณ(๐ฑ)(\psi _1^๐ณ(๐ฑ),\psi _2^๐ณ(๐ฑ))`$ is
$$\mathrm{\Psi }_f(๐ฑ)U\mathrm{\Psi }_m(๐ฑ)=\left(\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta \\ \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right)\mathrm{\Psi }_m(๐ฑ).$$
(1)
Each mass wave function is defined as a four-component spinorial function $`\psi _n(๐ฑ,t)`$, $`n=1,2`$ that satisfy the free Dirac equation
$$i\frac{}{t}\psi _n(๐ฑ,t)=H_n^D\psi _n(๐ฑ,t),n=1,2,$$
(2)
where the free Hamiltonian is the usual
$$H_n^Di๐ถ.+\beta m_n,n=1,2.$$
(3)
We will work in the flavor diagonal basis. This choice defines the flavor basis vectors simply as
$$\widehat{\nu }_e^๐ณ=(1,0),\widehat{\nu }_\mu ^๐ณ=(0,1),$$
(4)
while the flavor projectors are obviously
$$๐ฏ_{\nu _\alpha }\widehat{\nu }_\alpha \widehat{\nu }_\alpha ^๐ณ.$$
(5)
Actually, as an abuse of notation, the equivalence $`UU\text{1}_D`$ is implicit , as well as, $`๐ฏ_{\nu _\alpha }๐ฏ_{\nu _\alpha }\text{1}_D`$; the symbol $`\text{1}_D`$ refers to the identity matrix in spinorial space.
The total Hamiltonian governing the dynamics of $`\mathrm{\Psi }_m`$ is $`H^D=\mathrm{diag}(H_1^D,H_2^D)`$. From the considerations above, $`\mathrm{\Psi }_f(๐ฑ,t)`$ satisfy the equation
$$i\frac{}{t}\mathrm{\Psi }_f(๐ฑ,t)UH^DU^1\mathrm{\Psi }_f(๐ฑ,t).$$
(6)
The solution to the equation above can be written in terms of a flavor evolution operator $`K^D`$ as
$$\mathrm{\Psi }_f(๐ฑ,t)=K^D(t)\mathrm{\Psi }_f(๐ฑ,0)=d^3๐ฑ^{}K^D(๐ฑ๐ฑ^{};t)\mathrm{\Psi }_f(๐ฑ^{},0),$$
(7)
where
$$K^D(๐ฑ๐ฑ^{};t)=\frac{d^3๐ฉ}{(2\pi )^3}K^D(๐ฉ;t)e^{i๐ฉ.(๐ฑ๐ฑ^{})}.$$
(8)
We can calculate $`K^D(t)`$ in any representation (momentum or position) as
$`K^D(t)`$ $`=`$ $`Ue^{iH^Dt}U^1`$ (9)
$`=`$ $`\left(\begin{array}{cc}\mathrm{cos}^2\theta e^{iH_1^Dt}+\mathrm{sin}^2\theta e^{iH_2^Dt}& \mathrm{cos}\theta \mathrm{sin}\theta (e^{iH_1^Dt}e^{iH_2^Dt})\\ \mathrm{cos}\theta \mathrm{sin}\theta (e^{iH_1^Dt}e^{iH_2^Dt})& \mathrm{sin}^2\theta e^{iH_1^Dt}+\mathrm{cos}^2\theta e^{iH_2^Dt}\end{array}\right).`$ (12)
The conversion probability is then
$`๐ซ(\nu _e\nu _\mu ;t)`$ $`=`$ $`{\displaystyle ๐๐ฑ\mathrm{\Psi }_f^{}(๐ฑ,0)K^D(t)๐ฏ_{\nu _\mu }K^D(t)\mathrm{\Psi }_f(๐ฑ,0)}`$ (13)
$`=`$ $`{\displaystyle ๐๐ฉ\stackrel{~}{\psi }_{\nu _e}^{}(๐ฉ)(K_{\mu e}^D)^{}K_{\mu e}^D(๐ฉ,t)\stackrel{~}{\psi }_{\nu _e}(๐ฉ)},`$ (14)
satisfying the initial condition $`\mathrm{\Psi }_f^๐ณ(๐ฑ,0)=(\psi _{\nu _e}^๐ณ(๐ฑ,0),0)`$. Such initial condition implies, in terms of mass eigenfunctions, $`\psi _1(๐ฑ,0)=\psi _2(๐ฑ,0)=\psi _{\nu _e}(๐ฑ)`$, as a requirement to obtain an initial wave function with definite flavor ccnishi04 . The function $`\stackrel{~}{\psi }_{\nu _e}(๐ฉ)`$ denotes the inverse Fourier transform of $`\psi _{\nu _e}(๐ฑ)`$ (see Eqs. (110) and (111)).
Before obtaining the conversion probability for Dirac fermions, let us replace the spinorial functions $`\psi _n(๐ฑ)`$ by spinless one-component wave functions $`\phi _n(๐ฑ)`$ in the flavor wave function $`\mathrm{\Psi }_f^๐ณ(๐ฑ)(\phi _{\nu _e}(๐ฑ),\phi _{\nu _\mu }(๐ฑ))`$ and mass wave function $`\mathrm{\Psi }_m^๐ณ(๐ฑ)(\phi _1(๐ฑ),\phi _2(๐ฑ))`$. We also replace the Dirac Hamiltonian in momentum space $`H_n^D(๐ฉ)`$ (3) by the relativistic energy $`E_n(๐ฉ)=\sqrt{๐ฉ^2+m_n^2}`$. Inserting these replacements into Eq. (13) we can recover the usual oscillation probability ccnishi04 ; Bernardini:PRD05
$`๐ซ(\nu _e\nu _\mu ;t)`$ $`=`$ $`{\displaystyle ๐๐ฑ|\widehat{\nu }_\mu ^๐ณ\mathrm{\Psi }_f(๐ฑ,t)|^2}`$ (15)
$`=`$ $`{\displaystyle ๐๐ฉ|K_{\mu e}^S(๐ฉ,t)\stackrel{~}{\phi }_{\nu _e}(๐ฉ)|^2}`$ (16)
$`=`$ $`{\displaystyle ๐๐ฉ๐ซ(๐ฉ,t)|\stackrel{~}{\phi }_{\nu _e}(๐ฉ)|^2},`$ (17)
where $`\mathrm{\Psi }_f(๐ฑ,0)^๐ณ=(\phi _{\nu _e}(๐ฑ)^๐ณ,0)`$, $`K_{\mu e}^S(๐ฉ,t)(K^S)_{21}=\mathrm{sin}\theta \mathrm{cos}\theta (e^{iE_1(๐ฉ)t}e^{iE_2(๐ฉ)t})`$ and
$$๐ซ(๐ฉ,t)=\mathrm{sin}^22\theta \mathrm{sin}^2(\mathrm{\Delta }E\left(๐ฉ\right)t/2)$$
(18)
is just the standard oscillation formula. The conversion probability (15) in this case is then the standard oscillation probability smeared out by the initial momentum distribution. If the substitution $`|\stackrel{~}{\phi }_{\nu _e}(๐ฉ)|^2\delta ^3(๐ฉ๐ฉ_0)`$ is made the standard oscillation formula is recovered: it corresponds to the plane-wave limit.
After we have checked the standard oscillation formula can be recovered for spinless particles restricted to positive energies in the plane-wave limit, we can return to the case of Dirac fermions. We can obtain explicitly the terms of the mixed evolution kernel (9) by using the property of the Dirac Hamiltonian in momentum space $`H_{n}^{D}{}_{}{}^{2}=E_n^2(๐ฉ)\text{1}_D`$, which leads
$`(K_{\mu e}^D)^{}K_{\mu e}^D(๐ฉ,t)`$ $`=`$ $`๐ซ(๐ฉ,t)[1f(๐ฉ)]\text{1}_D`$ (20)
$`+\mathrm{sin}^22\theta f(๐ฉ)\mathrm{sin}^2(\overline{E}t)\text{1}_D,`$
where
$$f(๐ฉ)=\frac{1}{2}[1\frac{๐ฉ^2+m_1m_2}{E_1E_2}],$$
(21)
and $`๐ซ(๐ฉ,t)`$ is the standard conversion probability function (18). A unique implication of Eq. (20), which is proportional to the identity matrix in spinorial space, is that the conversion probability (13) does not depend on the spinorial structure of the initial flavor wave function but only on its momentum density as
$`๐ซ(\nu _e\nu _\mu ;t)={\displaystyle }d๐ฉ\{๐ซ(๐ฉ,t)[1f(๐ฉ)]`$
$`+\mathrm{sin}^22\theta f(๐ฉ)\mathrm{sin}^2(\overline{E}t)\}\stackrel{~}{\psi }^{}_{\nu _e}(๐ฉ)\stackrel{~}{\psi }_{\nu _e}(๐ฉ).`$ (22)
(The tilde will denote the inverse Fourier transformed function throughout this paper.) Furthermore, the modifications in Eq. (II) compared to the scalar conversion probability (15) are exactly the same modifications found in Ref. Bernardini:Euro and Ref. Bernardini:PRD05 after smearing out through a specific gaussian wave packet.
The conservation of total probability
$$๐ซ(\nu _e\nu _\mu ;t)+๐ซ(\nu _e\nu _e;t)=1,$$
(23)
is automatic in virtue of
$$K_{ee}^D(๐ฉ,t)K_{ee}^D(๐ฉ,t)+K_{\mu e}^D(๐ฉ,t)K_{\mu e}^D(๐ฉ,t)=\text{1}_D,$$
(24)
and the survival and conversion probability for an initial muon neutrino are identical to the probabilities for an initial electron neutrino because of the relations
$`K_{\mu \mu }^D(๐ฉ,t)K_{\mu \mu }^D(๐ฉ,t)`$ $`=`$ $`K_{ee}^D(๐ฉ,t)K_{ee}^D(๐ฉ,t),`$ (25)
$`K_{\mu e}^D(๐ฉ,t)K_{\mu e}^D(๐ฉ,t)`$ $`=`$ $`K_{e\mu }^D(๐ฉ,t)K_{e\mu }^D(๐ฉ,t).`$ (26)
To explain the origin of the additional oscillatory terms in Eq. (II) it is instructive to rewrite the free Dirac time evolution operator, in momentum space, in the form
$$e^{iH_n^Dt}=e^{iE_nt}\mathrm{\Lambda }_{n+}^D+e^{iE_nt}\mathrm{\Lambda }_n^D,$$
(27)
where
$$\mathrm{\Lambda }_{n\pm }^D=\frac{1}{2}(\text{1}_D\pm \frac{H_n^D}{E_n}),$$
(28)
are the projector operators to positive (+) or negative (-) energy eigenstates of $`H_n^D`$. By using the decomposition above (27), we can analyze $`K_{\mu e}^D`$ in Eq. (9), which contains the terms
$`e^{iH_1^Dt}e^{iH_2^Dt}`$ $`=`$ $`e^{i\mathrm{\Delta }Et}\mathrm{\Lambda }_{1+}^D\mathrm{\Lambda }_{2+}^D+e^{i\mathrm{\Delta }Et}\mathrm{\Lambda }_1^D\mathrm{\Lambda }_2^D`$ (30)
$`+e^{i2\overline{E}t}\mathrm{\Lambda }_{1+}^D\mathrm{\Lambda }_2^D+e^{i2\overline{E}t}\mathrm{\Lambda }_1^D\mathrm{\Lambda }_{2+}^D,`$
plus its hermitian conjugate. Since $`\mathrm{\Lambda }_{1\pm }^D\mathrm{\Lambda }_2^D0`$, it can be seen that the rapid oscillating terms come from the interference between, e.g., the positive frequencies of the Hamiltonian $`H_1^D`$ and negative energies of the Hamiltonian $`H_2^D`$. One may think that by restricting the initial wave function to contain only positive energy contributions would eliminate the rapid oscillatory terms, as zitterbewegung disappears for superpositions of solely positive energy states in Dirac theory IZ , but it does not happen. The positive energy eigenfunctions with respect to a basis characterized by a mass $`m_1`$ necessarily have non-null components of negative energy with respect to another basis characterized by $`m_2`$ (this point is illustrated in appendix B). Thus the rapid oscillatory contributions are an inevitable consequence of this framework and it is always present independently of the initial WP, if initially a definite flavor is chosen. However, its influence, quantified by the function $`f(๐ฉ)`$ in Eq. (21), is negligible for momentum distributions around ultra-relativistic values Bernardini:Euro . This rapid oscillatory terms will also be found for charged spin 0 particle oscillations in the next section, with contributions slightly different from the ones obtained for spin 1/2 particles.
### II.1 Inclusion of Left-Handedness
Until this point, we have been considering the oscillation of general flavor โparticle numberโ for general Dirac neutrinos. However, due to the left handed nature of weak interactions only left-handed components are produced and detected. To incorporate this fact into, for example, the conversion probability in Eq. (13), it is sufficient to use initial left-handed WPs and replace the kernel of Eq. (20) by the projected counterpart
$`LK_{\mu e}^D(๐ฉ,t)LK_{\mu e}^D(๐ฉ,t)L=๐ซ^D(๐ฉ,t)L`$
$`{\displaystyle \frac{1}{4}}\mathrm{sin}^22\theta \left({\displaystyle \frac{m_1}{E_1}}\mathrm{sin}E_1t{\displaystyle \frac{m_2}{E_2}}\mathrm{sin}E_2t\right)^2L,`$ (31)
where $`๐ซ^D(๐ฉ,t)=K_{\mu e}^D(๐ฉ,t)K_{\mu e}^D(๐ฉ,t)`$ is the conversion kernel of Eq. (20) and $`L=(1\gamma _5)/2`$ is the projector to left chirality. The conservation of total probability (23) no longer holds because there is a probability loss due to the undetected right handed component
$$LK_{\mu e}^DRK_{\mu e}^D(๐ฉ,t)L=\frac{1}{4}\mathrm{sin}^22\theta \left(\frac{m_1}{E_1}\mathrm{sin}E_1t\frac{m_2}{E_2}\mathrm{sin}E_2t\right)^2L,$$
(32)
where $`R=(1+\gamma _5)/2`$ is the projector to right chirality. We can see that the probability loss (32) is proportional to the ratio $`m_n^2/E_n^2`$ which is negligible for ultra-relativistic neutrinos. The total probability loss for an initial left-handed electron neutrino turning into right-handed neutrinos, irrespective of the final flavor, is given by the kernel
$$LK_{\mu e}^DRK_{\mu e}^D(๐ฉ,t)L+LK_{ee}^DRK_{ee}^D(๐ฉ,t)L=\left[\mathrm{cos}^2\theta \left(\frac{m_1}{E_1}\right)^2\mathrm{sin}^2E_1t+\mathrm{sin}^2\theta \left(\frac{m_2}{E_2}\right)^2\mathrm{sin}^2E_2t\right]L.$$
(33)
To obtain the unphysical complementary kernels responsible for the conversion of right-handed component to right-handed and left-handed components, it is enough to make the substitution $`LR`$ in all formulas.
## III Flavor Oscillation for Spin 0
The derivation of the usual conversion probability (15) takes into account only the positive frequency contributions. The mass wave function used to obtain Eq. (15) corresponds to the solutions of the wave equation
$$i\frac{}{t}\phi (๐ฑ,t)=\sqrt{^2+m^2}\phi (๐ฑ,t),$$
(34)
which is equivalent to the Dirac equation in the Foldy-Wouthuysen representation FW , restricted to positive energies. The evolution kernel for this equation is not satisfactory from the point of view of causality (Thaller, , p.18), i.e, the kernel is not null for spacelike intervals. Moreover, the eigenfunctions restricted to one sign of energy do not form a complete set FV .
To recover a causal propagation in the spin 0 case, the Klein-Gordon wave equation must be considered. In the first quantized version, the spectrum of the solutions have positive and negative energy as in the Dirac case. However, to take advantage of the Hamiltonian formalism used so far, it is more convenient do work in the Sakata-Taketani (ST) Hamiltonian formalism FV where each mass wave function is formed by two components
$$\mathrm{\Phi }_n(๐ฑ,t)=\left(\begin{array}{c}\phi _n(๐ฑ,t)\\ \chi _n(๐ฑ,t)\end{array}\right),n=1,2.$$
(35)
The components $`\phi `$ and $`\chi `$ are combinations of the usual scalar Klein-Gordon wave function $`\varphi (x)`$ and its time derivative $`_0\varphi (x)`$. This is necessary since the Klein-Gordon equation is a second order differential equation in time and the knowledge of the function and its time derivative is necessary to completely define the time evolution.
The time evolution in this formalism is governed by the Hamiltonian FV
$$H_n^{ST}=(\tau _3+i\tau _2)\frac{^2}{2m_n}+m_n^2,$$
(36)
which satisfies the condition $`(H_n^{ST})^2=(^2+m_n^2)\text{1}_{ST}`$, like the Dirac Hamiltonian (3). The $`\tau _k`$ represents the usual Pauli matrices and $`\text{1}_{ST}`$ is the identity matrix.
A charge density endnote1 can be defined as
$$\overline{\mathrm{\Phi }}_n\mathrm{\Phi }_n\mathrm{\Phi }_n^{}\tau _3\mathrm{\Phi }_n=|\phi _n|^2|\chi _n|^2,$$
(37)
which is equivalent to the one found in Klein-Gordon notation $`i\varphi ^{}\stackrel{}{}_0\varphi `$. Needless to say, this density (37) is only non-null for complex (charged) wave functions. The charge density $`\overline{\mathrm{\Phi }}\mathrm{\Phi }`$ is the equivalent of fermion probability density $`\psi ^{}\psi `$ in the Dirac case, although the former is not positive definite as the latter. The adjoint $`\overline{\mathrm{\Phi }}=\mathrm{\Phi }^{}\tau _3`$ were defined to make explicit the (non positive definite) norm structure of the conserved charge
$$๐๐ฑ\overline{\mathrm{\Phi }}_n(๐ฑ,t)\mathrm{\Phi }_n(๐ฑ,t)(\mathrm{\Phi }_n,\mathrm{\Phi }_n)=\text{time independent}.$$
(38)
Consequently, the adjoint of any operator $`\mathrm{\Omega }`$ can be defined as $`\overline{\mathrm{\Omega }}=\tau _3\mathrm{\Omega }^{}\tau _3`$, satisfying $`(\overline{\mathrm{\Omega }}\mathrm{\Phi },\mathrm{\Phi })=(\mathrm{\Phi },\mathrm{\Omega }\mathrm{\Phi })`$. Within this notation, the Hamiltonians of Eq. (36) is self-adjoint, $`\overline{H}_n^{ST}=H_n^{ST}`$, and the time invariance of Eq. (38) is assured.
We can assemble, as in the previous section, the mass wave functions into $`\mathrm{\Psi }_m^๐ณ(\mathrm{\Phi }_1^๐ณ,\mathrm{\Phi }_2^๐ณ)`$ and the flavor wave functions into $`\mathrm{\Psi }_f^๐ณ(\mathrm{\Phi }_{\nu _e}^๐ณ,\mathrm{\Phi }_{\nu _\mu }^๐ณ)`$, satisfying the mixing relation $`\mathrm{\Psi }_fU\mathrm{\Psi }_m`$. The equivalence of $`UU\text{1}_{ST}`$ and of $`๐ฏ_{\nu _\alpha }๐ฏ_{\nu _\alpha }\text{1}_{ST}`$ are implicit without modification in the notations. Then, the time evolution of $`\mathrm{\Psi }_f`$ can be given through a time evolution operator $`K^{ST}`$ acting in the same form as in Eq. (7). In complete analogy to the calculations from Eq. (8) to Eq. (13), we can define the conversion probability as
$`๐ซ(\nu _e\nu _\mu ;t)`$ $`=`$ $`{\displaystyle ๐๐ฑ\overline{\mathrm{\Psi }}_f(๐ฑ,0)\overline{K^{ST}(t)}๐ฏ_{\nu _\mu }K^{ST}(t)\mathrm{\Psi }_f(๐ฑ,0)}`$ (39)
$`=`$ $`{\displaystyle ๐๐ฉ\overline{\stackrel{~}{\mathrm{\Phi }}_e(๐ฉ)}\overline{K_{\mu e}^{ST}}K_{\mu e}^{ST}(๐ฉ,t)\stackrel{~}{\mathrm{\Phi }}_e(๐ฉ)},`$ (40)
where $`\mathrm{\Psi }_f(๐ฑ,0)^๐ณ=(\mathrm{\Phi }_e(๐ฑ)^๐ณ,0)`$. The adjoint operation were also extended to $`\overline{\mathrm{\Psi }}_f=\mathrm{\Psi }_f^{}(\text{1}_\theta \tau _3)`$, where $`\text{1}_\theta `$ is the identity in mixing space.
The information of time evolution, hence oscillation, is all encoded in
$`\overline{K_{\mu e}^{ST}}K_{\mu e}^{ST}(๐ฉ,t)`$ $`=`$ $`๐ซ(๐ฉ,t)[1f(\mu ๐ฉ)]\text{1}_{ST}`$ (42)
$`+\mathrm{sin}^22\theta f(\mu ๐ฉ)\mathrm{sin}^2(\overline{E}t)\text{1}_{ST},`$
where the function $`f(๐ฉ)`$ were already defined in Eq. (21) and
$$\mu =\sqrt{\frac{1}{2}(\frac{m_1}{m_2}+\frac{m_2}{m_1})}.$$
(43)
The factor $`\mu 1`$ determines the difference with the Dirac case in Eq. (20). The equality $`\mu =1`$ holds when $`m_1=m_2`$, i.e., when there is no oscillation.
## IV Connection with quantum field theory
The main improvement of the covariant approaches developed in secs. II and III is that the propagation kernels governed by Dirac and Sakata-Taketani Hamiltonians are causal, i.e., are null for spacelike separations (see Eqs. (127) and (128) and Refs. Thaller ; IZ ; Roman ). On the contrary, the kernel of spinless particles restricted only to positive energies is not null for spacelike intervals Thaller . From the point of view of relativistic classical field theories, a causal kernel guarantees, by the Cauchy theorem, the causal connection between the wave-function in two spacelike surfaces at different times Roman .
To compare the IWP and EWP approaches it is useful to rewrite the Dirac evolution kernel for a fermion of mass $`m_n`$, present in Eq. (7), in the form (IZ, , p.89)
$`K_n^D(xy)`$ $`=`$ $`\underset{๐ }{}\frac{d^3p}{2E_n}[u_n^s(x;๐ฉ)\overline{u}_n^s(y;๐ฉ)+v_n^s(x;๐ฉ)\overline{v}_n^s(y;๐ฉ)]\gamma _0`$ (44)
$``$ $`iS(xy;m_n)\gamma _0,n=1,2,`$ (45)
where $`(xy)^0=t,(xy)^i=(๐ฑ๐ฑ^{})^i`$ when compared to the notation of Eq. (7). The spinorial functions $`u,v`$, are the free solutions of the Dirac equation and they are explicitly defined in appendix A. (More familiar forms for the function $`S`$ are also shown in appendix A.) Clearly the function $`iS(xy;m_n)=0|\{\psi _n(x),\overline{\psi }_n(y)\}|0`$ satisfies the homogeneous Dirac equation with mass $`m_n`$ (2) and it is known to be null for spacelike intervals $`(xy)^2<0`$ Thaller ; Roman .
In contrast, the Feynman propagator $`iS_F(xy)`$ appears in QFT. It is a Green function for the inhomogeneous Dirac equation obeying particular boundary conditions. The EWP approaches use this Green function for the propagation of virtual neutrinos. To directly compare the Feynman propagator to the kernel in Eq. (44) we can write $`iS_F`$ in the form
$`iS_F(xy;m_n)`$ $``$ $`0|T(\psi _n(x),\overline{\psi }_n(y))|0`$ (46)
$`=`$ $`\underset{๐ }{}\frac{d^3p}{2E_n}[u_n^s(x;๐ฉ)\overline{u}_n^s(y;๐ฉ)\theta (x_0y_0)`$ (48)
$`v_n^s(x;๐ฉ)\overline{v}_n^s(y;๐ฉ)\theta (y_0x_0)].`$
Although the function $`S_F`$ is called causal propagator, it is not null for spacelike intervals, and it naturally arises in QFT when interactions are present and treated in a covariant fashion. Equation (46) shows that the propagator $`S_F`$ describes positive energy states propagating forward in time and negative energy states propagating backward in time (IZ, , p.91). At a first glance, both neutrino and antineutrino parts of Eq. (46) seem to contribute to the space-time integrations present in covariant perturbation theory, as neutrino-antineutrino contributions in Eq. (44) have led to Eq. (II).
In the following we will show in an EWP approach that for large separations between production and detection both neutrino and antineutrino parts may contribute as intermediate neutrinos for certain situations.
We will follow the calculations made in Ref. Dolgov , using, instead of the scalar interaction, the effective charged-current weak lagrangian
$`_W`$ $`=`$ $`G{\displaystyle \underset{i,\alpha =1}{\overset{N=3}{}}}[\overline{l}_\alpha (x)\gamma ^\mu LU_{\alpha i}\nu _i(x)J_\mu (x)`$ (50)
$`+\overline{\nu }_i(x)U_{\alpha i}^{}\gamma ^\mu Ll_\alpha (x)J_\mu ^{}(x)]`$
$`=`$ $`_1+_1^{},`$ (51)
where $`G=\sqrt{2}G_F`$ and $`J_\mu `$ is the sum of any effective leptonic or hadronic current. The lagrangian (50) is written only in terms of physical mass eigenstate fields, which coincides with flavor eigenstate fields only for the charged leptons: $`l_1(x)e(x),l_2(x)\mu (x),\mathrm{}`$ .
Suppose the process Dolgov ; Giunti:qft:02 where a charged lepton $`l_\alpha `$ hit a nucleus A turning it into another nucleus A with emission of a neutrino (this process happens around $`x_A`$). Subsequently the neutrino travels a long distance and hit a nucleus B which transforms into B emitting a lepton $`l_\beta `$ (this process happens around $`x_B`$). The whole process looks like $`l_\alpha +A+Bl_\beta +A^{}+B^{}`$ with transition amplitude given by
$$A^{}(๐ฉ_A^{}),B^{}(๐ฉ_B^{}),l_\beta (๐ฉ_\beta )|S|A,B,l_\alpha .$$
(52)
The final states are momentum eigenstates while the initial states are localized Dolgov . The lowest order nonzero contribution of the scattering matrix $`S`$ to Eq. (52) is second order in the lagrangian (50). More explicitly, the term that contributes to the amplitude (52) comes from
$`S^{(2)}`$ $`=`$ $`{\displaystyle \frac{i^2}{2}}T_W^2={\displaystyle \frac{1}{2}}T_1+_1^{}^2`$ (53)
$``$ $`T_1_1^{}`$ (54)
$``$ $`G^2{\displaystyle d^4xd^4y\underset{\beta \alpha }{}_{\beta \alpha }(x,y)},`$ (55)
where $``$ stands for space-time integration and
$$_{\beta \alpha }(x,y)\underset{i}{}:J_\mu (x)\overline{l}_\beta (x)\gamma ^\mu LU_{\beta i}iS_F(xy;m_i)U_{\alpha i}^{}\gamma ^\nu Ll_\alpha (y)J_\nu ^{}(y):.$$
(56)
In Eq. (54) we kept only the mixed product and in Eq. (55) we kept from all possible terms in Wick expansion (IZ, , $`\mathrm{p}.180`$) only the term responsible for the transition of interest.
Then the transition amplitude (52) can be calculated as
$`G^2A^{}(๐ฉ_A^{}),`$ $`B^{}(๐ฉ_B^{}),l_\beta (๐ฉ_\beta )|S^{(2)}|A,B,l_\alpha `$ (57)
$`={\displaystyle d^4yd^4xB^{}(๐ฉ_B^{})|J_\mu (y)|BA^{}(๐ฉ_A^{})|J_\nu ^{}(x)|A}`$ (58)
$`\times \overline{u}_\beta (y,๐ฉ_\beta )\gamma ^\mu L{\displaystyle \underset{i}{}}U_{\beta i}U_{\alpha i}^{}iS_F(yx;m_i)\gamma ^\nu L0|l_\alpha (x)|l_\alpha `$ (59)
$`{\displaystyle \underset{i}{}}U_{\beta i}U_{\alpha i}^{}๐_i.`$ (60)
The initial states must be chosen in such a way that $`A,l_\alpha `$ are localized around $`x_A=(t_A,๐ฑ_A)`$ and $`B`$ is localized around $`x_B=(t_B,๐ฑ_B)`$, since we are ultimately interested in large separations $`|๐ฑ_B๐ฑ_A|`$. We can implement explicitly those localization conditions into the wave packets
$`B^{}(๐ฉ_B^{})|J_\mu (y)|B`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^{\frac{3}{2}}}}{\displaystyle \widehat{d๐ช_B}e^{ip_B^{}.y}J_\mu ^{BB^{}}(๐ช_B,๐ฉ_B^{})\psi _B(๐ช_B)e^{iq_B.(yx_B)}}`$ (61)
$`A^{}(๐ฉ_A^{})|J_\nu ^{}(x)|A`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^{\frac{3}{2}}}}{\displaystyle \widehat{d๐ช_A}e^{ip_A^{}.y}J_\mu ^{AA^{}}(๐ช_A,๐ฉ_A^{})\psi _A(๐ช_A)e^{iq_A.(xx_A)}}`$ (62)
$`0|l_\alpha (x)|l_\alpha `$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^{\frac{3}{2}}}}{\displaystyle \widehat{d๐ช_\alpha }\psi _\alpha (๐ช_\alpha )e^{iq_\alpha .(xx_A)}},`$ (63)
where $`\widehat{d๐ช}=d๐ช(2E(๐ช))^{1/2}`$, $`J_\mu ^{BB^{}}(๐ช_B,๐ฉ_B^{})=B^{}(๐ฉ_B^{})|J_\mu (0)|B(๐ช_B)`$ and $`J_\nu ^{AA^{}}(๐ช_A,๐ฉ_A^{})=A^{}(๐ฉ_A^{})|J_\nu ^{}(0)|A(๐ช_A)`$.
Following the calculations from Eq. (59) with the localization aspects of Eqs. (61)-(63) included, we arrive at
$`๐_i`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^6}}{\displaystyle \widehat{d๐ช_B}\widehat{d๐ช_A}\widehat{d๐ช_\alpha }J_\mu ^{BB^{}}(๐ช_B,๐ฉ_B^{})\psi _B(๐ช_B)J_\mu ^{AA^{}}(๐ช_A,๐ฉ_A^{})\psi _A(๐ช_A)}`$ (65)
$`e^{iq_B.x_B}e^{i(q_A+q_\alpha ).x_A}\overline{u}_\beta (๐ฉ_\beta )\gamma ^\mu L\left[{\displaystyle d^4xd^4ye^{i\kappa _\beta .y}e^{i\kappa _\alpha .x}iS_i(yx)}\right]\gamma ^\nu L\psi _\alpha (๐ช_\alpha ),`$
where $`\kappa _\beta =(\kappa _\beta ^0,๐ฟ_\beta )`$, $`\kappa _\alpha =(\kappa _\alpha ^0,๐ฟ_\alpha )`$ and
$`๐ฟ_\beta `$ $`๐ฉ_\beta +๐ฉ_B^{}๐ช_B,\kappa _\beta ^0E_\beta (๐ฉ_\beta )+E_B^{}(๐ฉ_B^{})E_B(๐ช_B),`$ (66)
$`๐ฟ_\alpha `$ $`๐ฉ_\alpha ๐ฉ_A^{}+๐ช_A,\kappa _\alpha ^0E_\alpha (๐ช_\alpha )E_A^{}(๐ฉ_A^{})+E_A(๐ช_A).`$ (67)
By using the results of Eqs. (152) and (153) the expression between square brackets in Eq. (65) gives
$`2\pi \delta (\kappa _\beta ^0\kappa _\alpha ^0){\displaystyle }`$ $`d๐ฑd๐ฒ{\displaystyle \frac{i}{4\pi r}}e^{ik_\omega r}e^{i๐ฟ_\beta .๐ฒ}e^{i๐ฟ_\alpha .๐ฑ}`$ (68)
$`\times [u_i(k_\omega \widehat{๐ซ})\overline{u}_i(k_\omega \widehat{๐ซ})\theta (\omega _im_i)v_i(k_\omega \widehat{๐ซ})\overline{v}_i(k_\omega \widehat{๐ซ})\theta (\omega _im_i)],`$ (69)
where $`r|๐ฒ๐ฑ|`$, $`\widehat{๐ซ}(๐ฒ๐ฑ)/r`$, $`\omega _i\kappa _\beta ^0=\kappa _\alpha ^0`$ and $`k_\omega \sqrt{\omega _i^2m_i^2}`$. The crucial point here is that, depending on the masses and momenta of the incoming particles, both neutrinos ($`u\overline{u}`$) and antineutrinos ($`v\overline{v}`$) can contribute to the amplitude (65) depending on the sign of its energy $`\omega _i`$, restricted to $`|\omega _i|>m_i`$; the off-shell contributions for $`\omega _i[m_i,m_i]`$ are exponentially decreasing and then negligible for large distances (see appendix C). We will see in the following that antineutrino contributions in this case is possible and it corresponds to unphysical contributions.
We are interested in large productionโdetection separations. It permits us to approximate, as in Ref. Dolgov , $`rR+\widehat{๐}.(๐ฒ๐ฑ_B)\widehat{๐}.(๐ฑ๐ฑ_A)`$ and $`\widehat{๐ซ}\widehat{๐}`$, where $`R|๐ฑ_B๐ฑ_A|`$ and $`\widehat{๐}(๐ฑ_B๐ฑ_A)/R`$. Such approximations inserted in Eq. (69) lead to momentum conservation on $`x_A`$ and $`x_B`$ vertices:
$`2\pi \delta (\kappa _\beta ^0\kappa _\alpha ^0)`$ $`{\displaystyle \frac{i}{4\pi R}}e^{ik_\omega R}e^{ik_\omega \widehat{๐}.(๐ฑ_B๐ฑ_A)}(2\pi )^3\delta ^3(๐ฟ_\beta k_\omega \widehat{๐})(2\pi )^3\delta ^3(๐ฟ_\alpha k_\omega \widehat{๐})`$ (70)
$`\times [u_i(k_\omega \widehat{๐})\overline{u}_i(k_\omega \widehat{๐})\theta (\omega _im_i)v_i(k_\omega \widehat{๐})\overline{v}_i(k_\omega \widehat{๐})\theta (\omega _im_i)].`$ (71)
At this point we have all the information to analyze whether the antineutrino part of the propagator contributes to the overall process. Neither of the isolated processes $`A+l_\alpha A^{}+\overline{\nu }_i`$ and $`B+\overline{\nu }_iB^{}+l_\beta `$ are allowed if we calculate the transition amplitude for them separately using the weak Lagrangian (50). (For Majorana neutrinos they are strongly suppressed by helicity mismatch.) So far four-momentum conservation in both $`x_A`$ and $`x_B`$ vertices were automatically required from the calculations; among them the requirement of energy conservation for intermediate neutrinos with respect to the accompanying particles in vertex $`x_A`$ ($`\omega _i=\kappa _\alpha ^0`$) and in vertex $`x_B`$ ($`\omega _i=\kappa _\beta ^0`$), is already implicit. The remaining are explicit in the delta functions of Eq. (71). The on-shell condition for neutrinos ($`|\omega _i|^2k_\omega ^2=m_i^2`$) for long distance propagation was also automatic. What the calculations did not required is a definite sign for $`\omega _i`$, for all possible momenta constrained by the mentioned energy-momentum conservations. To analyze if and under what conditions both signs are possible is equivalent to study the kinematics of two-body to two-body scattering allowing the sign of one particle energy to be free. Putting in equations, for vertex $`x_A`$, assuming the particle $`A`$ at rest, we obtain from $`(p_Ap_i)^2=(p_A^{}p_\alpha )^2`$ the neutrino energy
$$E_i=\frac{1}{2M_A}[M_A^2M_A^{}^2+m_i^2m_\alpha ^2+2E_\alpha E_A^{}2๐ฉ_\alpha .๐ฉ_A^{}].$$
(72)
The minimum value of right-hand side of Eq. (72) corresponds to last two terms equal to $`2m_\alpha M_A^{}`$, which gives for the minimum
$$\mathrm{min}(E_i)=\frac{1}{2M_A}[M_A^2(M_A^{}m_\alpha )^2+m_i^2].$$
(73)
The values $`\omega _i=\kappa _\alpha ^0`$ are bounded from below by the value in Eq. (73). Imposing $`\mathrm{min}(E_i)>m_i`$ and $`\mathrm{min}(E_i)<m_i`$ is respectively equivalent to
$`M_AM_A^{}`$ $`>`$ $`(m_\alpha m_i)`$ (74)
$`M_AM_A^{}`$ $`<`$ $`(m_\alpha +m_i),`$ (75)
for $`M_A>m_i`$ and $`M_A^{}>m_\alpha `$. It is clear that depending on the value of the masses, condition (75) may be satisfied leading to antineutrino contributions to Eq. (65) for a range of possible incoming momenta. Of course the condition (74) is sufficient to exclude antineutrino contributions but it also excludes the cases where a threshold energy is required for the lepton $`l_\alpha `$ to initiate the production reaction. Thus to prevent antineutrino contributions, it is better to adopt the weaker condition of restricting the sign of the energy of intermediate neutrinos $`\omega _i`$ to be positive, keeping only the first term in Eq. (69). Analogous analysis lead to possible momenta and mass values that allow $`\kappa _\beta ^0<m_i`$ for vertex $`x_B`$, still compatible with $`\kappa _\alpha ^0=\kappa _\beta ^0`$. Notice that condition $`\kappa _\alpha ^0>m_i`$ is exactly the kinematical condition to allow the production of physical neutrinos in $`x_A`$ and $`\kappa _\beta ^0>m_i`$ allow only the contribution of neutrinos with energy above threshold to trigger the detection reaction. The violation of these conditions implies in kinematically impossible contributions in production or detection.
Restricted to condition $`\omega _i>0`$ we can insert the expression above into Eq. (65) which yields
$`๐_i`$ $`=`$ $`{\displaystyle \widehat{d๐ช_\alpha }2\pi \delta (\kappa _\beta ^0\kappa _\alpha ^0)\theta (\omega _im_i)\frac{i}{4\pi R}e^{ik_\omega Ri\omega _i(t_Bt_A)}e^{i(p_k+p_B^{}).x_B}e^{ip_A^{}.x_A}}`$ (78)
$`\times u_k(๐ฉ_\beta )\gamma ^\mu Lu_i(k_\omega \widehat{๐})\overline{u}_i(k_\omega \widehat{๐})\gamma ^\nu L\psi _\alpha (๐ช_\alpha )`$
รJฮผBB(๐ชB,๐ฉB)ฯB(๐ชB)EB(๐ชB)JฮฝAA(๐ชA,๐ฉA)ฯA(๐ชA)EA(๐ชA)|
=qB-+pฮฒpBkฯ^R
=qA+-pAqฮฑkฯ^R
.absentevaluated-atsubscriptsuperscript๐ฝ๐ตsuperscript๐ต๐subscript๐ช๐ตsubscriptsuperscript๐ฉ๐ตsubscript๐๐ตsubscript๐ช๐ตsubscript๐ธ๐ตsubscript๐ช๐ตsubscriptsuperscript๐ฝ๐ดsuperscript๐ด๐subscript๐ช๐ดsubscriptsuperscript๐ฉ๐ดsubscript๐๐ดsubscript๐ช๐ดsubscript๐ธ๐ดsubscript๐ช๐ด
=qB-+pฮฒpBkฯ^R
=qA+-pAqฮฑkฯ^R
\displaystyle\times\left.J^{BB^{\prime}}_{\mu}({\bf q}_{B},{\bf p}^{\prime}_{B})\frac{\psi_{B}({\bf q}_{B})}{\sqrt{E_{B}({\bf q}_{B})}}J^{AA^{\prime}}_{\nu}({\bf q}_{A},{\bf p}^{\prime}_{A})\frac{\psi_{A}({\bf q}_{A})}{\sqrt{E_{A}({\bf q}_{A})}}\right|_{\parbox{54.24994pt}{
$\mbox{\scriptsize${\bf q}_{B}={\bf p}_{\beta}+{\bf p}^{\prime}_{B}-k_{\omega}\hat{{\bf R}}$}$
\\
$\mbox{\scriptsize${\bf q}_{A}={\bf p}^{\prime}_{A}-{\bf q}_{\alpha}+k_{\omega}\hat{{\bf R}}$}$
}}\quad.
Notice that the step function $`\theta (\omega _im_i)`$ prevents non-physical neutrinos to contribute to the process.
Particularly, if we use a unidimensional wave packet for the incoming lepton $`l_\alpha `$
$$\psi _\alpha (๐ช)=\psi _\alpha (q_x,q_y,q_z)=\delta (q_x)\delta (q_y)\psi _{\alpha z}(q_z),$$
(79)
we obtain an amplitude analogous to Ref. Dolgov
$`{\displaystyle \underset{i}{}}U_{\beta i}U_{\alpha i}^{}๐_i`$ $`=`$ $`{\displaystyle \underset{i}{}}{\displaystyle \frac{i}{4\pi R}}e^{ik_\omega Ri\omega _i(t_Bt_A)}2\pi \left|2\mathrm{p}_\alpha {\displaystyle \frac{}{๐ช_\alpha ^2}}(\kappa _\beta ^0\kappa _\alpha ^0)\right|_{๐ช_\alpha =\mathrm{p}_\alpha \widehat{z}}^1e^{i(p_k+p_B^{}).x_B}e^{ip_A^{}.x_A}`$ (82)
$`\times U_{\beta i}u_\beta (๐ฉ_\beta )\gamma ^\mu Lu_i(k_\omega \widehat{๐})U_{\alpha i}^{}\overline{u}_i(k_\omega \widehat{๐})\gamma ^\nu L\psi _{\alpha z}(\mathrm{p}_\alpha )`$
รJฮผBB(๐ชB,๐ฉB)ฯB(๐ชB)EB(๐ชB)JฮฝAA(๐ชA,๐ฉA)ฯA(๐ชA)EA(๐ชA)|
=qB-+pฮฒpBkฯ^R
=qA+-pApฮฑ^zkฯ^R
,absentevaluated-atsubscriptsuperscript๐ฝ๐ตsuperscript๐ต๐subscript๐ช๐ตsubscriptsuperscript๐ฉ๐ตsubscript๐๐ตsubscript๐ช๐ตsubscript๐ธ๐ตsubscript๐ช๐ตsubscriptsuperscript๐ฝ๐ดsuperscript๐ด๐subscript๐ช๐ดsubscriptsuperscript๐ฉ๐ดsubscript๐๐ดsubscript๐ช๐ดsubscript๐ธ๐ดsubscript๐ช๐ด
=qB-+pฮฒpBkฯ^R
=qA+-pApฮฑ^zkฯ^R
\displaystyle\times\left.J^{BB^{\prime}}_{\mu}({\bf q}_{B},{\bf p}^{\prime}_{B})\frac{\psi_{B}({\bf q}_{B})}{\sqrt{E_{B}({\bf q}_{B})}}J^{AA^{\prime}}_{\nu}({\bf q}_{A},{\bf p}^{\prime}_{A})\frac{\psi_{A}({\bf q}_{A})}{\sqrt{E_{A}({\bf q}_{A})}}\right|_{\parbox{57.26382pt}{
$\mbox{\scriptsize${\bf q}_{B}={\bf p}_{\beta}+{\bf p}^{\prime}_{B}-k_{\omega}\hat{{\bf R}}$}$
\\
$\mbox{\scriptsize${\bf q}_{A}={\bf p}^{\prime}_{A}-\mathrm{p}_{\alpha}\hat{z}+k_{\omega}\hat{{\bf R}}$}$
}}~{},~{}~{}
where $`\mathrm{p}_\alpha `$ is the root of $`f(|๐ช_\alpha |=\mathrm{p}_\alpha )=\kappa _\beta ^0\kappa _\alpha ^0=0`$, which comes from energy conservation from the whole process; if there is no root the process is kinematically forbidden. The detection probability is proportional to the square of the amplitude (82) integrated over the final phase space $`d๐ฉ_A^{}d๐ฉ_B^{}d๐ฉ_\beta [2E_A^{}(๐ฉ_A^{})2E_B^{}(๐ฉ_B^{})2E_\beta (๐ฉ_\beta )]^1`$. In particular, since $`๐ฉ_\beta ,๐ฉ_A^{},๐ฉ_B^{}`$ are fixed, the phases that differ for different intermediate neutrinos $`\nu _i`$ are only $`k_\omega R\omega _i(t_At_B)`$ which is the same result obtained in Ref. Dolgov, (except for terms which depend on the mean velocity of particles $`A`$ and $`B`$).
So far we have shown in an EWP approach both processes in $`x_A`$ and $`x_B`$ should be considered real scattering processes with real neutrinos involved. The off-shell contributions are negligible to large distances and antineutrino contributions were explicitly excluded by eliminating the second term of Eq. (69). These informations permit us to rewrite Eq. (78) in a slightly different form
$`G^2{\displaystyle \underset{i}{}}U_{\beta i}U_{\alpha i}^{}๐_i`$ $`=`$ $`{\displaystyle \underset{i}{}}{\displaystyle \frac{d๐ฉ}{2E_i(๐ฉ)}d^4yB^{}(๐ฉ_B^{}),l_\beta (๐ฉ_\beta )|_1(y)e^{i(Pp_i).x_B}|B,\nu _i(๐ฉ)}`$ (84)
$`{\displaystyle d^4x\theta (yx)A^{}(๐ฉ_A^{}),\nu _i(๐ฉ)|_1^{}(x)e^{iP.x_A}|A,l_\alpha },`$
where $`P=(H,๐)`$ is the energy-momentum operator. A change of notation were made here: in Eq. (82) the states $`|B`$ and $`|A,l_\alpha `$ are centered around the origin while in Eqs. (52)-(63) they are respectively centered around $`x_B`$ and $`x_A`$; the translation is explicitly performed by the translation operator $`e^{iP.x}`$. Additionally, the step function $`\theta (yx)`$ is necessary to ensure that the contributions of points $`y`$ around $`x_B`$ should always be after the contributions of points $`x`$ around $`x_A`$. By following the same steps from Eq. (59) to Eq. (79) we can arrive from Eq. (84) to Eq. (82).
Equation (84) shows us the amplitude of the entire process from production to detection in โdecomposedโ form (apart from the step function in time): the amplitude of production process multiplied by the amplitude of detection process summed over all possible intermediate real neutrinos of different masses $`m_i`$ and momentum $`๐ฉ`$. (The sum over spins are implicit.)
### IV.1 A simple second quantized formulation
Considering that only real neutrinos or antineutrinos (one of them exclusively) should travel from production to detection, the possibility to use the free second quantized theory for spin 1/2 fermions to describe flavor oscillations is investigated in this section. This simple, type B and QFT description of flavor oscillation phenomena guarantees only particle or antiparticle propagation, keeping the nice property of giving normalized oscillation probabilities, like the first quantized examples treated in secs. II and III.
To accomplish the task of calculating oscillation probabilities in QFT we have to define the neutrino states that are produced and detected through weak interactions. Firstly, we define the shorthand for the combination of fields appearing in the weak effective charged-current lagrangian (50)
$$\nu _\alpha (x)U_{\alpha i}\nu _i(x),\alpha =e,\mu .$$
(85)
We will restrict the problem to two flavor families and use the matrix $`U`$ as the same in Eq. (1). The mass eigenfields $`\nu _i(x)`$, $`i=1,2`$, are the physical fields for which the mass eigenstates $`|\nu _i(๐ฉ)`$ are well defined asymptotic states. The free fields $`\nu _i(x)`$ can be expanded in terms of creation and annihilation operators (see appendix A) and the projection to the one-particle space defines the mass wave function
$$\psi _{\nu _i}(๐ฑ;g_i)=0|\nu _i(๐ฑ)|\nu _i:g_i\underset{๐ }{}๐๐ฉ\frac{g_i^s(๐ฉ)}{\sqrt{2E_i}}u_i^s(๐ฑ;๐ฉ),i=1,2,$$
(86)
where
$$|\nu _i:g_i\underset{๐ }{}๐๐ฉg_i^s(๐ฉ)|\nu _i(๐ฉ,s).$$
(87)
Since the creation operators for neutrinos (antineutrinos) can be written in terms of the free fields $`\overline{\nu }_i(x)`$ ($`\nu _i(x)`$), we can define the โflavorโ states as the superpositions of mass eigenstates
$`|\nu _\alpha :\{g\}`$ $``$ $`U_{\alpha i}^{}|\nu _i:g_i`$ (88)
$`|\overline{\nu }_\alpha :\{g\}`$ $``$ $`U_{\alpha i}|\overline{\nu }_i:g_i.`$ (89)
The details of creation are encoded in the functions $`g_i`$.
We can also define
$`\psi _{\nu _\alpha \nu _e}(x;\{g\})`$ $``$ $`0|\nu _e(x)|\nu _\alpha :\{g\}`$ (90)
$`=`$ $`U_{ei}U_{\alpha i}^{}\psi _{\nu _i}(x;g_i),`$ (91)
where $`\psi _{\nu _i}(x)`$ are then mass wave functions defined in Eq. (86). We can see from Eq. (90) that if $`\psi _{\nu _1}(๐ฑ,t)=\psi _{\nu _2}(๐ฑ,t)=\psi (๐ฑ)`$, for a given time $`t`$, $`\psi _{\nu _e\nu _e}(๐ฑ,t)=\psi (๐ฑ)`$ and $`\psi _{\nu _\mu \nu _e}(๐ฑ,t)=0`$ due to the unitarity of the mixing matrix.
Although this approach does not rely on flavor Fock spaces and Bogoliubov transformations, we can use the same observables used by Blasone and Vitiello to quantify flavor oscillation BV:obs : the flavor charges, which are defined as
$$Q_\alpha (t)=d๐ฑ:\nu _\alpha ^{}(๐ฑ,t)\nu _\alpha (๐ฑ,t):,\alpha =e,\mu ,$$
(92)
where $`::`$ denotes normal ordering. Note that the $`Q_e(t)+Q_\mu (t)=Q`$ is conserved BV:AP95 , the two flavor charges are compatible for equal times, i.e., $`[Q_e(t),Q_\mu (t)]=0`$, and $`\nu :\{g\}|Q|\nu :\{g\}=\pm \nu :\{g\}|\nu :\{g\}`$ for any particle state (+) or antiparticle state (-). Notice that in the second quantized version the charges can acquire negative values, despite the fermion probability density in first quantization is a positive definite quantity. The conservation of total charge guarantees the conservation of total probability (23).
We can further split the flavor charges into left-handed (-) and right-handed (+) parts
$$Q_\alpha ^{(\pm )}(t)=d๐ฑ:\nu _\alpha ^{}(๐ฑ,t)\frac{1}{2}(\text{1}\pm \gamma _5)\nu _\alpha (๐ฑ,t):,\alpha =e,\mu ,$$
(93)
where $`Q_\alpha ^{(+)}+Q_\alpha ^{()}=Q_\alpha `$. These components will be used to calculate the left-handed to right-handed transition.
With the flavor charges defined, we can calculate, for example, the conversion probability
$`๐ซ(\nu _e\nu _\mu ;t)`$ $``$ $`\nu _e:\{g\}|Q_\mu (t)|\nu _e:\{g\}`$ (94)
$`=`$ $`U_{\mu i}U_{\mu j}^{}U_{ej}U_{ei}^{}{\displaystyle ๐๐ฉe^{i(E_iE_j)t}\stackrel{~}{\psi }_{\nu _j}^{}(๐ฉ;g_j)\stackrel{~}{\psi }_{\nu _i}(๐ฉ;g_i)},`$ (95)
where the neutrino wave functions $`\psi _{\nu _i}`$ are defined in terms of the function $`g_i(๐ฉ)`$ in Eq. (86). If we could equate the two mass wavefunctions in momentum space $`\stackrel{~}{\psi }_{\nu _1}(๐ฉ;g_1)=\stackrel{~}{\psi }_{\nu _2}(๐ฉ;g_2)=\stackrel{~}{\psi }_{\nu _e}(๐ฉ)`$ we would obtain, from Eq. (95), the standard two family conversion probability (15)
$$๐ซ(\nu _e\nu _\mu ;t)=๐๐ฉ๐ซ(๐ฉ,t)\stackrel{~}{\psi }_{\nu _e}^{}(๐ฉ)\stackrel{~}{\psi }_{\nu _e}(๐ฉ),$$
(96)
where $`๐ซ`$ was defined in Eq. (18). However, the equality can not hold as proved in appendix (B): two wavefunctions with only positive energy components with respect to two bases characterized by different masses can not be equal. Then, it is not possible to impose a flavor definite condition. Instead, we can write
$$g_i(๐ฉ,s)=\frac{u_i^s(๐ฉ)}{\sqrt{2E_i(๐ฉ)}}\stackrel{~}{\psi }_i(๐ฉ),$$
(97)
where $`\stackrel{~}{\psi _i(๐ฉ)}`$ is the initial wave function associated to the neutrino of mass $`m_i`$ at creation, taking care to maintain the normalization $`๐๐ฉ|g_i(๐ฉ)|^2=1`$; any transition amplitude can be written in the form Eq. (97). In general $`\stackrel{~}{\psi }_i(๐ฉ)=\stackrel{~}{\psi }(๐ฉ,m_i)`$, and then, for small mass differences,
$$\stackrel{~}{\psi }_i(๐ฉ)\stackrel{~}{\psi }(๐ฉ,\overline{m})\pm \frac{\mathrm{\Delta }m}{2}\frac{}{\overline{m}}\stackrel{~}{\psi }(๐ฉ,\overline{m}),$$
(98)
where $`\overline{m}=(m_1+m_2)/2`$ and $`\mathrm{\Delta }m=m_2m_1`$. Keeping only the first term, $`\stackrel{~}{\psi }(๐ฉ,\overline{m})\stackrel{~}{\psi }(๐ฉ)`$, we obtain from Eq. (95),
$`๐ซ(\nu _e\nu _\mu ;t)`$ $`=`$ $`{\displaystyle ๐๐ฉ๐ซ(๐ฉ,t)\stackrel{~}{\psi }^{}(๐ฉ)[\text{1}\frac{1}{2}\mathrm{\Lambda }_1(๐ฉ)\frac{1}{2}\mathrm{\Lambda }_2(๐ฉ)]\stackrel{~}{\psi }(๐ฉ)}`$ (100)
$`+\frac{1}{4}\mathrm{sin}^22\theta {\displaystyle }d๐ฉ\stackrel{~}{\psi }^{}(๐ฉ)[f(๐ฉ)\mathrm{cos}(\mathrm{\Delta }Et)i{\displaystyle \frac{\mathrm{\Delta }m}{2E_1E_2}}๐ธ.๐ฉ\mathrm{sin}(\mathrm{\Delta }Et)]\stackrel{~}{\psi }(๐ฉ).`$
Notice that in this case, the conversion probability is non-null for $`t=0`$,
$$๐ซ(\nu _e\nu _\mu ;0)=\frac{1}{4}\mathrm{sin}^22\theta ๐๐ฉf(๐ฉ)\stackrel{~}{\psi }^{}(๐ฉ)\stackrel{~}{\psi }(๐ฉ),$$
(101)
which imply a direct lepton flavor violation in creation. However, since $`f(๐ฉ)(\mathrm{\Delta }m)^2/(4๐ฉ^2)`$ for ultra-relativistic momenta, the violation is hopelessly feeble for direct measurement. Among the deviations of the conversion probability (100) compared to the standard one (96), only the last term is of order $`\mathrm{\Delta }m/\overline{E}`$, the rest is of order $`(\mathrm{\Delta }m/\overline{E})^2`$ (the contributions of $`\mathrm{\Lambda }_i`$ can be estimated by $`[v_2^{}(๐ฉ,s)u_1(๐ฉ,s^{})]^2๐ฉ^2[\mathrm{\Delta }m+\mathrm{\Delta }E]^2/[(m_1+E_1)(m_1+E_1)]`$). Even so, $`\mathrm{\Delta }m/\overline{E}10^9`$ for $`\mathrm{\Delta }m^210^3\mathrm{eV}^2`$, $`\overline{m}1/2\mathrm{e}\mathrm{V}`$ and $`\overline{E}1\mathrm{M}\mathrm{e}\mathrm{V}`$, which can not be seen in actual oscillation experiments. Nevertheless, it is important to note that the knowledge of $`\mathrm{\Delta }m`$ in conjunction with $`\mathrm{\Delta }m^2`$ gives information about the absolute mass scale because of $`\mathrm{\Delta }m^2=2\overline{m}\mathrm{\Delta }m`$.
Using $`Q_\alpha ^{()}`$ of Eq. (93) instead of $`Q_\alpha `$ in Eq. (95) and $`\stackrel{~}{\psi }(๐ฉ)=L\stackrel{~}{\psi }(๐ฉ)`$ in Eq. (96) we obtain
$$๐ซ(\nu _{eL}\nu _{\mu R};t)=๐๐ฉ\left[\frac{m_1m_2}{4E_1E_2}๐ซ(๐ฉ,t)+\frac{1}{4}\mathrm{sin}^22\theta \left(\frac{m_1}{2E_1}\frac{m_2}{2E_2}\right)^2\right]\stackrel{~}{\psi }^{}(๐ฉ)\stackrel{~}{\psi }(๐ฉ).$$
(102)
The total probability loss from the conversion of initial left-handed electron neutrino to right-handed neutrinos yields
$$๐ซ(\nu _{eL}\nu _{eR};t)+๐ซ(\nu _{eL}\nu _{\mu R};t)=๐๐ฉ\left[\mathrm{cos}^2\theta \left(\frac{m_1}{2E_1}\right)^2+\mathrm{sin}^2\theta \left(\frac{m_2}{2E_2}\right)^2\right]\stackrel{~}{\psi }^{}(๐ฉ)\stackrel{~}{\psi }(๐ฉ).$$
(103)
Notice Eq. (103) does not depend on time in contrast to its first quantized analog in Eq. (33). Other conversion and survival probabilities can be obtained from Eq. (24) and
$$๐ซ(\nu _{eL}\nu _{\mu R};t)+๐ซ(\nu _{eL}\nu _{\mu L};t)=๐ซ(\nu _e\nu _\mu ;t).$$
(104)
The exchange of $`LR`$ does not modify the formulas, provided we also change the chirality of the initial wave function.
For completeness we calculate the additional conversion probabilities including the second term of Eq. (98)
$`\delta ๐ซ(\nu _e\nu _\mu ;t)`$ $`=`$ $`\frac{1}{4}\mathrm{sin}^22\theta {\displaystyle \frac{\mathrm{\Delta }m}{2}}{\displaystyle }d๐ฉ\stackrel{~}{\psi }^{}(๐ฉ)[{\displaystyle \frac{H_2}{2E_2}}{\displaystyle \frac{H_1}{2E_1}}+{\displaystyle \frac{\mathrm{\Delta }m}{2E_1E_2}}๐ธ.๐ฉ`$ (106)
$`+(\mathrm{\Lambda }_{1+}+\mathrm{\Lambda }_{2+}f(๐ฉ))i\mathrm{sin}\mathrm{\Delta }Et]{\displaystyle \frac{}{\overline{m}}}\stackrel{~}{\psi }(๐ฉ)+h.c.,`$
$`\delta ๐ซ(\nu _{eL}\nu _{\mu R};t)`$ $`=`$ $`\frac{1}{4}\mathrm{sin}^22\theta {\displaystyle \frac{\mathrm{\Delta }m}{2}}{\displaystyle }d๐ฉ\stackrel{~}{\psi }^{}(๐ฉ)[\left({\displaystyle \frac{m_2}{2E_2}}\right)^2\left({\displaystyle \frac{m_1}{2E_1}}\right)^2`$ (108)
$`+{\displaystyle \frac{m_1m_2}{2E_1E_2}}i\mathrm{sin}\mathrm{\Delta }Et]{\displaystyle \frac{}{\overline{m}}}\stackrel{~}{\psi }(๐ฉ)+h.c.,`$
which have terms of order $`\mathrm{\Delta }m`$ and $`(\mathrm{\Delta }m)^2`$.
To calculate the conversion probability for antineutrinos $`\overline{\nu }_e\overline{\nu }_\mu `$, it is enough to use
$$g_i^s(๐ฉ)\stackrel{~}{\psi }_i^{}(๐ฉ)\frac{v_i^s(๐ฉ)}{\sqrt{2E_i(๐ฉ)}},$$
(109)
instead of Eq. (97), replace $`tt`$ and $`\stackrel{~}{\psi }(๐ฉ)\frac{๐ถ.๐ฉ}{|๐ฉ|}\stackrel{~}{\psi }(๐ฉ)`$ in the expressions corresponding to neutrinos (95)โ(108), or apply charge conjugation $`\stackrel{~}{\psi }(๐ฉ)\gamma _0C\stackrel{~}{\psi }^{}(๐ฉ)`$. These prescriptions can be inferred from direct comparison to the calculations and beware that the definition of antineutrino states are defined with $`g_i^s(๐ฉ)`$ (139).
The formulas obtained in this second quantized version does not have the interference terms between positive and negative energies like in Eq. (II). Such interference terms are absent because the possible mixed terms like $`b_2(๐ฉ)a_1^{}(๐ฉ)|0`$ are null for an initial โflavorโ state superposition that contains only particle states (or only antiparticles states). Notice that the irrelevance of the initial spinorial structure no longer holds in this second quantized version, which can be seen, for example, in Eq. (100).
## V Discussion and Conclusions
Using the Dirac equation which is a relativistic covariant equation we obtained oscillation probabilities respecting causal propagation. The oscillation formulas obtained had additional rapid oscillating terms depending on the frequency $`E_1+E_2`$, with respect to the usual oscillation formulas with wave packets. Such additional oscillatorial character seemed to have its origin in the intrinsic spinorial character used. However, we have seen that such terms also appear in the charged spin 0 particle oscillations. In fact the rapid oscillatorial terms arise from the interference of positive and negative frequency parts of the initial WP and they are always present independently of the initial wave packet if initially a flavor definite condition were imposed. In addition, within Dirac theory, we have shown the detailed spinorial character of the initial wave function was irrelevant for flavor oscillation, independently of the momenta involved. The inclusion of the left-handed nature of the created and detected neutrinos could also be simply achieved. It is important to stress that the modifications found in this context would have tiny observable effects to the flavor oscillation of ultrarelativistic neutrinos.
Regarding second quantized approaches (sec. IV), in particular, EWP approaches, we can compare the propagators used in the latter to the evolution kernels in IWP approaches. Both the free evolution kernel and the Feynman propagator for fermions contain the contribution from particle and antiparticles. From this perspective, EWP approaches could also contain both contributions from neutrinos and antineutrinos, as in the first quantized approaches presented in secs. II and III. To analyze this point, an EWP calculation were carried out explicitly in sec. IV following Ref. Dolgov, . Restricted to a case where only neutrinos would be present, the calculation showed that the antineutrino contribution were not excluded automatically in the formalism but a subsidiary condition could be necessary: the sign of the frequency of the intermediate neutrinos should be restricted to be positive. In such case, there can be interference terms between positive and negative frequencies, possibly yielding rapid oscillation terms similar to the ones obtained in Dirac theory of sec. II. However, it should be stressed that the origin of the interference between positive and negative terms are different in the first quantized Dirac theory treated in sec. II and in the EWP (second quantized) treated in this sec. IV. The former comes from causality and completeness arguments in a classical field theory perspective, while the latter have its origin in the consideration of non-physical contributions. The restriction implied by the subsidiary condition automatically guarantee that: (i) only real neutrinos that are kinematically allowed in production contributes and (ii) in detection, only neutrinos with energies above threshold to trigger the detection reaction contribute. Otherwise, kinematically forbidden reactions in production or detection could be possible through exchange of virtual antineutrinos carrying negative energy. The overall energy-momentum conservation, though, is always respected (smeared out through finite momentum distributions) through production/propagation/detection processes. Since the presence of both neutrinoโantineutrino contributions is common to all EWP approaches, the subsidiary condition necessary in the EWP approach analyzed is possibly necessary in any approach with virtual neutrino propagation. (Unless a stronger condition like Eq. (74) is already implicit.) For example, in Eq. (14) of Ref. GS, , the subsidiary condition (for antineutrinos) is satisfied because the detection reaction is an elastic scattering. (Although the detection electrons are off-shell (bound state), the subsidiary condition is valid.) An important remark in this respect is that the calculations of the production and detection amplitudes as separate processes take automatically into account the physical kinematical conditions (i) and (ii) through the energy-momentum delta functions. It is also important to stress that the result obtained here did not depend on particular wave packets neither on the particular interaction used. The interesting point is that by imposing such subsidiary condition beforehand enables us to write the transition amplitude for the entire production/propagation/detection process in a โdecomposedโ form, with simple physical interpretation. A more detailed discussion about the decomposition of the process in separate production, propagation and detection processes can be found in Ref. Cardall.00, . Similar conclusions can be drawn for the case where only antineutrinos should propagate: the sign to be chosen should be negative. A realistic EWP approach for antineutrino propagation is given in Ref. GS, . To conclude this part, EWP approaches for neutrino oscillations require for consistency, but do not automatically imply, real intermediate neutrinos or antineutrinos exclusively.
All the properties discussed above about EWP approaches suggest that the description of two macroscopically distant scattering processes (production and detection) as a single scattering process described by a single scattering matrix have to be treated with care. We can be confident about the use of the $`S`$ matrix to describe any microscopic event through perturbation theory to any order of expansion (if the theory is renormalizable), but the extension to macroscopically distant reactions is not automatic. Actually, if the two processes are indeed not causally connected it can be proved that the $`S`$ matrix decomposes as the product of the two $`S`$ matrices for the two distant and independent processes Weinberg .
From the considerations above, the positive and negative interference terms in the first quantized approaches considered seem unphysical. To support this idea, it was shown in sec. IV.1 that a simple second quantized, type B, and IWP-like, approach could be devised using the second quantized free theory maintaining the simple properties of IWP approaches but eliminating the undesirable interference of positive and negative frequencies that was inevitable in the relativistic quantum mechanical context. On the other hand, new ingredients such as the direct flavor violation in creation and deviations from the standard oscillation formula were found. The deviations include the probability loss due to the conversion of left-handed neutrinos to right-handed neutrinos. Unfortunately, those new effects are tiny either because of the small mass difference or the ultrarelativistic nature of neutrinos and they are not feasible for direct observation in actual oscillation experiments.
###### Acknowledgements.
This work was supported by Conselho Nacional de Desenvolvimento Cientรญfico e Tecnolรณgico (CNPq). The author would like to thank Prof. J. C. Montero for a critical reading of the manuscript.
## Appendix A Notation and definitions
The (scalar, spinorial or ST) wave functions related by Fourier transforms are denoted as
$`\phi (๐ฑ)`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^{3/2}}}{\displaystyle ๐๐ฉ\stackrel{~}{\phi }(๐ฉ)e^{i๐ฉ.๐ฑ}},`$ (110)
$`\stackrel{~}{\phi }(๐ฉ)`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^{3/2}}}{\displaystyle ๐๐ฑ\phi (๐ฑ)e^{i๐ฉ.๐ฑ}}.`$ (111)
The tilde denotes the inverse Fourier transformed function.
Using the property of the Dirac or ST Hamiltonian, $`H_n^2=(๐ฉ^2+m_n)^2\text{1}`$, we can write the evolution operator in the form
$$e^{iH_nt}=\mathrm{cos}(E_nt)i\frac{H_n}{E_n}\mathrm{sin}(E_nt),$$
(112)
where the momentum dependence have to be replaced by $`i`$ in coordinate space.
The free neutrino field expansion used is ($`i=1,2`$)
$$\nu _i(x)=\underset{๐ }{}\frac{d๐ฉ}{2E_๐ฉ}[u_i^s(x;๐ฉ)a_i(๐ฉ,s)+v_i^s(x;๐ฉ)b_i^{}(๐ฉ,s)],$$
(113)
where the creation and annihilation operators satisfy the canonical anticommutation relations
$`\{a_i(๐ฉ,r),a_j^{}(๐ฉ^{},s)\}`$ $`=`$ $`\delta _{ij}\delta _{rs}2E_i(๐ฉ)\delta ^3(๐ฉ๐ฉ^{}),`$ (114)
$`\{b_i(๐ฉ,r),b_j^{}(๐ฉ^{},s)\}`$ $`=`$ $`\delta _{ij}\delta _{rs}2E_i(๐ฉ)\delta ^3(๐ฉ๐ฉ^{});`$ (115)
all other anticommutation relations are null. The functions $`u,v`$ are defined as
$`u_i^s(x;๐ฉ)`$ $`=`$ $`u_i^s(๐ฉ){\displaystyle \frac{e^{ip_i.x}}{(2\pi )^{3/2}}},`$ (116)
$`u_i^s(๐ฉ)`$ $`=`$ $`{\displaystyle \frac{m_i+E_i\gamma ^0๐ฉ.๐ธ}{\sqrt{E_i+m_i}}}u_0^s,`$ (117)
$`v_i^s(x;๐ฉ)`$ $`=`$ $`v_i^s(๐ฉ){\displaystyle \frac{e^{ip_i.x}}{(2\pi )^{3/2}}},`$ (118)
$`v_i^s(๐ฉ)`$ $`=`$ $`{\displaystyle \frac{m_iE_i\gamma ^0+๐ฉ.๐ธ}{\sqrt{E_i+m_i}}}v_0^s,`$ (119)
where $`p_i.x=E_i\left(๐ฉ\right)t๐ฉ.๐ฑ`$ and they satisfy the properties
$`\overline{u}_0^ru_0^s=u_0^ru_0^s`$ $`=`$ $`\overline{v}_0^rv_0^s=v_0^rv_0^s=\delta _{rs},`$ (120)
$`v_0^ru_0^s`$ $`=`$ $`u_0^rv_0^s=0r,s,`$ (121)
$`\underset{๐ }{}u_i^s(๐ฉ)\overline{u_i}^s(๐ฉ)`$ $`=`$ $`\mathit{}+m_i=2E_i(๐ฉ)\mathrm{\Lambda }_{i+}^D(๐ฉ)\gamma ^0,`$ (122)
$`\underset{๐ }{}v_i^s(๐ฉ)\overline{v_i}^s(๐ฉ)`$ $`=`$ $`\mathit{}m_i=2E_i(๐ฉ)\mathrm{\Lambda }_i^D(๐ฉ)\gamma ^0.`$ (123)
The Feynman propagator for fermions is
$`iS_F(xy)`$ $``$ $`0|T(\psi (x)\overline{\psi }(y))|0`$ (124)
$`=`$ $`{\displaystyle \frac{d^4p}{\left(2\pi \right)^4}\frac{i}{\mathit{}m+iฯต}e^{ip.(xy)}}`$ (125)
$`=`$ $`(i\partial ฬธ+m)i\mathrm{\Delta }_F(xy;m).`$ (126)
The function $`S`$ in Eq. (44) and its equivalent for the Sakata-Taketani Hamiltonian can be written as
$`iS(x;m)`$ $`=`$ $`(i\partial ฬธ+m)i\mathrm{\Delta }(x;m),`$ (127)
$`K^{ST}(x;m)`$ $`=`$ $`[i_t{\displaystyle \frac{^2}{2m}}(\tau _3+i\tau _2)+m^2]i\mathrm{\Delta }(x;m),`$ (128)
$`i\mathrm{\Delta }(x;m)`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^3}}{\displaystyle d^4p\delta (p^2m^2)ฯต(p_0)e^{ip.x}}`$ (129)
$`=`$ $`{\displaystyle \frac{1}{(2\pi )^3}}{\displaystyle \frac{d๐ฉ}{2E_p}[e^{ip.x}e^{+ip.x}]}.`$ (130)
The free neutrino eigenstates are defined as
$`|\nu _i(๐ฉ,s)`$ $``$ $`{\displaystyle \frac{a_i^{}(๐ฉ,s)}{\sqrt{2E_i}}}|0`$ (131)
$`=`$ $`{\displaystyle ๐๐ฑ\nu _i^{}(x)|0\frac{u_i(x;๐ฉ)}{\sqrt{2E_i}}},`$ (132)
$`|\overline{\nu }_i(๐ฉ,s)`$ $``$ $`{\displaystyle \frac{b_i^{}(๐ฉ,s)}{\sqrt{2E_i}}}|0`$ (133)
$`=`$ $`{\displaystyle ๐๐ฑ\frac{v_i^{}(x;๐ฉ)}{\sqrt{2E_i}}\nu _i(x)|0},`$ (134)
whose normalization is $`\nu _j(๐ฉ,r)|\nu _i(๐ฉ^{},s)=\delta _{ij}\delta _{rs}\delta ^3(๐ฉ๐ฉ^{})`$. The same normalization is valid for the antiparticle states. The states with finite momentum distributions are defined as
$`|\nu _i:g`$ $`=`$ $`\underset{๐ }{}๐๐ฉg^s(๐ฉ)|\nu _i(๐ฉ,s)`$ (135)
$`=`$ $`{\displaystyle ๐๐ฑ\nu _i^{}(x)|0\psi _{\nu _i}(x)},`$ (136)
$`\psi _{\nu _i}(x)`$ $``$ $`\underset{๐ }{}๐๐ฉg^s(๐ฉ)\frac{u_i^s(x;๐ฉ)}{\sqrt{2E_i}},`$ (137)
$`e^{iHt}|\nu _i:g`$ $`=`$ $`{\displaystyle ๐๐ฑ\nu _i^{}(๐ฑ,0)|0\psi _{\nu _i}(๐ฑ,t)},`$ (138)
$`|\overline{\nu }_i:g`$ $`=`$ $`\underset{๐ }{}๐๐ฉg^s(๐ฉ)|\overline{\nu }_i(๐ฉ,s)`$ (139)
$`=`$ $`{\displaystyle ๐๐ฑ\psi _{\overline{\nu }_i}^{}(x)\nu _i(x)|0},`$ (140)
$`\psi _{\overline{\nu }_i}(x)`$ $``$ $`\underset{๐ }{}๐๐ฉg^s(๐ฉ)\frac{v_i^s(x;๐ฉ)}{\sqrt{2E_i}},`$ (141)
$`e^{iHt}|\overline{\nu }_i:g`$ $`=`$ $`{\displaystyle ๐๐ฑ\psi _{\overline{\nu }_i}^{}(๐ฑ,t)\nu _i(๐ฑ,0)|0}.`$ (142)
## Appendix B Decomposition with respect to two bases
It is possible to decompose a given spinorial function $`\psi (๐ฑ)`$ in terms of bases depending on different masses $`m_1`$ and $`m_2`$. Equating
$$\psi (๐ฑ)=\frac{d๐ฉ}{2E_i}[u_i^s(๐ฑ;๐ฉ)g_i^{(+)}(๐ฉ,s)+v_i^s(๐ฑ;๐ฉ)g_i^{()}(๐ฉ,s)],i=1,2,$$
(143)
the expansion coefficients can be obtained
$`g_i^{(+)}(๐ฉ,s)`$ $`=`$ $`{\displaystyle ๐๐ฑu_i^s(๐ฑ;๐ฉ)\psi (๐ฑ)},`$ (144)
$`g_i^{()}(๐ฉ,s)`$ $`=`$ $`{\displaystyle ๐๐ฑv_i^s(๐ฑ;๐ฉ)\psi (๐ฑ)}.`$ (145)
From Eq. (145) we see that imposing the conditions
$`g_1^{()}(๐ฉ,s)`$ $`=`$ $`v_1^s(๐ฉ)\stackrel{~}{\psi }(๐ฉ)=0,`$ (146)
$`g_2^{()}(๐ฉ,s)`$ $`=`$ $`v_2^s(๐ฉ)\stackrel{~}{\psi }(๐ฉ)=0,`$ (147)
for all $`๐ฉ`$, lead to the equivalent conditions
$`[(m_1+E_2)(m_2+E_2)]v_0^s\stackrel{~}{\psi }(๐ฉ)`$ $`=`$ $`0,s=1,2,`$ (148)
$`\left[{\displaystyle \frac{1}{m_1+E_2}}{\displaystyle \frac{1}{m_2+E_2}}\right]v_0^s๐ธ.๐ฉ\stackrel{~}{\psi }(๐ฉ)`$ $`=`$ $`0,s=1,2,`$ (149)
where the property of Eq. (119) and $`\gamma _0v_0^s=v_0^s`$ were used. In case $`m_1m_2`$, we can use the decomposition $`\stackrel{~}{\psi }(๐ฉ)=\stackrel{~}{\psi }_+(๐ฉ)+\stackrel{~}{\psi }_{}(๐ฉ)`$, where $`\stackrel{~}{\psi }_\pm (๐ฉ)=(\text{1}\pm \gamma _0)\stackrel{~}{\psi }(๐ฉ)`$/2, and obtain from Eqs. (148) and (149) the conditions
$`v_0^s\stackrel{~}{\psi }_{}(๐ฉ)`$ $`=`$ $`0,s=1,2,`$ (150)
$`u_0^s๐.๐ฉ\stackrel{~}{\psi }_+(๐ฉ)`$ $`=`$ $`0,s=1,2,`$ (151)
where the properties $`๐ธ=\gamma _0\gamma _5๐`$ and $`u_0^s=\gamma _5v_0^s`$ were used in Eq. (151). The equations (150) and (151) are only satisfied if $`\stackrel{~}{\psi }_+(๐ฉ)=\stackrel{~}{\psi }_{}(๐ฉ)=0`$, since $`๐.๐ฉ`$ has only non-null eigenvalues and it commutes with $`\gamma _0`$. It is easier to reach this conclusion in the helicity basis $`\{u_0^{(\pm )},v_0^{(\pm )}\}`$ characterized by $`๐.๐ฉu_0^{(\pm )}=\pm |๐ฉ|u_0^{(\pm )}`$, but the result is basis independent.
## Appendix C Integrals
Splitting the Feynman propagator into positive and negative frequency parts $`iS_F(x)=iS^{(+)}(x)+iS^{()}(x)`$ the following integrals give us
$`{\displaystyle ๐te^{i\omega t}iS^{(+)}(๐ซ,t;m)}`$ $`=`$ $`(i){\displaystyle \frac{e^{ik_\omega r}}{4\pi r}}[(\omega \gamma _0k_\omega (\widehat{๐ซ}.๐ธ\mathbf{)}+m]\theta (\omega m)`$ (152)
$`{\displaystyle ๐te^{i\omega t}iS^{()}(๐ซ,t;m)}`$ $`=`$ $`(i){\displaystyle \frac{e^{ik_\omega r}}{4\pi r}}[(\omega \gamma _0k_\omega (\widehat{๐ซ}.๐ธ\mathbf{)}+m]\theta (\omega m),`$ (153)
where $`k_\omega =\sqrt{\omega ^2m^2},๐ซ=r\widehat{๐ซ}`$, the conditions $`mr,k_\omega r1`$ were assumed and terms behaving as $`1/r^2`$ were neglected.
To illustrate the calculations, Eq. (152) is obtained by
$`{\displaystyle ๐te^{i\omega t}iS^{(+)}(๐ซ,t;m)}`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^3}}{\displaystyle \frac{d๐ฉ}{2E(๐ฉ)}\frac{i(E(๐ฉ)\gamma _0๐ฉ.๐ธ+m)}{\omega E(๐ฉ)+iฯต}e^{i๐ฉ.๐ซ}}`$ (154)
$`=`$ $`{\displaystyle \frac{1}{(2\pi )^2}}{\displaystyle \frac{i}{2r}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}dp{\displaystyle \frac{p}{E(p)}}\{\mathrm{sin}pr{\displaystyle \frac{E(p)\gamma _0+m}{E(p)\omega iฯต}}`$ (156)
$`+[\mathrm{cos}pr{\displaystyle \frac{\mathrm{sin}pr}{pr}}]{\displaystyle \frac{ip(\widehat{๐ซ}.๐ธ)}{E(p)\omega iฯต}}\}.`$
In Eq. (154) the following identity is used
$$_0^{\mathrm{}}๐te^{\pm iEt}=\frac{\pm i}{E\pm iฯต}.$$
(157)
To get to the result of Eq. (152) it is necessary to split the functions $`\mathrm{sin}pr`$ and $`\mathrm{cos}pr`$ in Eq. (156) into exponentials and, for the $`e^{ipr}`$ part, integrate along a closed path formed by a half semicircle in the upper-half complex plane going round the branching line $`[im,i\mathrm{})`$ (for the $`e^{ipr}`$ part take the path reflected by the line defined by $`\mathrm{Re}z=0`$). The contribution from the paths beside the branching line yields a function which decreases more rapidly than $`e^{mr}`$ and it is negligible for $`mr1`$. The contributions for $`m<\omega <m`$ give a function with dependence $`e^{|k_\omega |r}`$ which is also negligible for large separations $`r`$.
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# Systematic Low-Energy Effective Theory for Magnons and Charge Carriers in an Antiferromagnet
## 1 Introduction
Almost 20 years after the discovery of high-temperature superconductivity in layered cuprates , identifying the dynamical mechanism behind it remains one of the great challenges in condensed matter physics. Ordinary low-temperature superconductors are weakly coupled electron systems in which phonon exchange mediates an attractive interaction that can overcome the Coulomb repulsion between electrons. As massless Nambu-Goldstone bosons of the spontaneously broken translation symmetry, phonons provide a natural mechanism for Cooper pair formation at low energies which is successfully quantified in BCS theory. In contrast to ordinary superconductors, layered high-$`T_c`$ cuprates are systems of strongly correlated electrons to which the weak coupling BCS theory is not readily applicable. Furthermore, the high transition temperatures of cuprate superconductors and the smallness of the isotope effect suggest that mechanisms other than phonon exchange may be responsible for Cooper pair formation. Since high-temperature superconductors are antiferromagnets before doping, it is natural to suspect (but not generally accepted) that magnons โ the Nambu-Goldstone bosons of the spontaneously broken $`SU(2)_s`$ spin symmetry โ may be important for binding electrons or holes into preformed pairs.
Even if spin fluctuations were not the key to explaining high-temperature superconductivity, the dynamics of charge carriers in an antiferromagnet is an interesting topic in itself. There is a vast literature on this subject. The dynamics of holes in an antiferromagnet has been investigated, for example, in . Understanding the dynamics of even just a single hole propagating in an antiferromagnet is a challenging problem. One can gain qualitative insight from a picture in which holes hop from site to site, leaving a string of flipped spins behind and thus locally destroying the antiferromagnetic order. Since the string costs energy proportional to its length, one might expect the holes to even be confined and thus have infinite mass. However, the locally destroyed antiferromagnetic order may be healed by appropriate hole hopping which renders the hole mass finite . Angle resolved photoemission spectroscopy experiments as well as a number of theoretical investigations indicate that the minimum of the dispersion (i.e. of the energy) of a single hole corresponds to lattice momenta $`(\pm \frac{\pi }{2},\pm \frac{\pi }{2})`$ in the Brillouin zone.
As one adds a second hole, the situation becomes more controversial. For example, there seems to be no consensus on the question if a pair of holes can form a bound state or not. If it can, the condensation of such pairs would provide a potential mechanism for high-temperature superconductivity. The effective theory to be constructed here can be used to analytically calculate the long-range magnon-mediated forces between holes using perturbation theory. It is very interesting to ask what happens when one dopes an antiferromagnet with a non-zero density of holes. At sufficient doping, experiments show that high-temperature superconductivity may arise. It has been argued on theoretical grounds that even an infinitesimal amount of doping may affect the antiferromagnetic phase and turn it into a spiral phase . A systematic investigation of this question is also possible using the effective theory of this paper, but it will require the use of non-perturbative methods.
The standard models for antiferromagnets and high-temperature superconductors are the Hubbard and $`t`$-$`J`$ model. Since these models are strongly coupled, they are not accessible to a systematic analytic treatment. As a consequence, analytic calculations in Hubbard-type models usually involve some uncontrolled approximations. Unfortunately, due to a severe fermion sign problem, away from half-filling these models can currently also not be simulated reliably. Hence, although they may indeed contain the relevant physics, Hubbard-type models have not yet led to a quantitative understanding of high-$`T_c`$ materials. An alternative to a microscopic description using Hubbard-type models is provided by phenomenological models formulated directly in terms of magnon and electron or hole fields . Although they may provide qualitative insight, such models do not lead to unambiguous predictions. In this paper, for the first time we introduce a systematic low-energy effective field theory for magnons and charge carriers in an antiferromagnet. Based only on symmetries and their spontaneous breakdown, the effective theory makes universal predictions for the entire class of antiferromagnetic cuprates. Although the effective theory is not renormalizable, it yields unambiguous results in a systematic low-energy expansion. In each order of the expansion, the results depend only on a finite number of material specific low-energy parameters whose values can be determined experimentally. The effective theory is not based on a specific microscopic model Hamiltonian but is universally applicable. Furthermore, and most important, in contrast to the strongly correlated electrons of Hubbard-type models, the electrons and holes of the effective field theory are quasi-particles that are weakly coupled to the magnons. Consequently, one may expect that the effective theory is more easily solvable than the underlying microscopic models.
Possible basic applications of the effective theory to be constructed in this paper include magnon-magnon, magnon-hole, and magnon-electron scattering as well as the determination of long-range magnon-mediated forces between the charge carriers. More ambitious applications could aim at a quantitative explanation of the Mott insulator state, the reduction of the staggered magnetization upon doping, the formation of a spiral phase, or at a systematic study of potential mechanisms for the preformation of electron or hole pairs in the antiferromagnetic phase. When such pairs condense they may become the Cooper pairs of high-temperature superconductivity. Except for a derivation of the dispersion relation of charge carriers, in this paper we do not consider applications yet, but concentrate entirely on the construction of the effective theory itself.
The construction in this paper is inspired by similar developments in the theory of the strong interactions. In contrast to the high-$`T_c`$ problem, where the choice of a microscopic model is controversial, there is general agreement that Quantum Chromodynamics (QCD) provides the correct microscopic description of the strong interactions. Still, similar to Hubbard-type models, solving QCD is notoriously hard. At โhalf-fillingโ, i.e. in the filled quark Dirac sea that represents the QCD vacuum, the $`SU(2)_LSU(2)_R`$ chiral symmetry of massless up and down quarks is spontaneously broken to the isospin symmetry $`SU(2)_{L=R}`$, resulting in three massless Nambu-Goldstone pions. This is analogous to the spontaneous breaking of the $`SU(2)_s`$ spin symmetry down to $`U(1)_s`$ that leads to antiferromagnetism. The corresponding Nambu-Goldstone bosons โ in this case two magnons โ are thus analogous to the pions of the strong interactions. It is possible to study chiral symmetry breaking in the QCD vacuum in numerical simulations of lattice QCD, just as it is possible to study antiferromagnetism by simulating the Hubbard model at half-filling. However, it is very useful to also investigate these phenomena with effective field theories. The low-energy effective theory for pions was pioneered by Weinberg and formulated as a systematic expansion in Gasserโs and Leutwylerโs chiral perturbation theory . Based on symmetry considerations and the observation that chiral symmetry is spontaneously broken, chiral perturbation theory makes rigorous predictions about the pion dynamics in terms of a few low-energy parameters such as the pion decay constant, the chiral condensate, and the Gasser-Leutwyler coefficients. Once these parameters are determined, either experimentally or through lattice QCD calculations, the effective theory makes unambiguous predictions in the low-energy domain.
Chiral perturbation theory can be applied to any Nambu-Goldstone phenomenon, and has indeed been used for both ferro- and antiferromagnetic magnons . To lowest order, for antiferromagnetic magnons the low-energy parameters of chiral perturbation theory are the spin stiffness $`\rho _s`$ and the spin-wave velocity $`c`$. At low energies chiral perturbation theory describes all aspects of the magnon dynamics just in terms of these two parameters. For example, the low-energy physics of the Hubbard model at half-filling is completely described by the effective theory once $`\rho _s`$ and $`c`$ have been determined in terms of the Hubbard model parameters $`t`$ and $`U`$.
A numerical challenge in high-$`T_c`$ physics is to simulate the Hubbard model away from half-filling. This requires a solution of the corresponding fermion sign problem. Similarly, simulating lattice QCD at non-zero baryon chemical potential, i.e. after โdopingโ the QCD vacuum with quarks, is prevented by a severe complex action problem. Like for high-$`T_c`$ materials at sufficient doping, one expects that QCD at sufficiently high baryon density becomes a superconductor, in that case for the color charge carried by quarks and gluons . In contrast to high-temperature superconductivity, the mechanism responsible for color-superconductivity is well understood in terms of one-gluon exchange. Color-superconductivity requires very large baryon densities and may thus arise only in the core of compact neutron or quark stars. However, superconductivity โ not of color but of ordinary electric charge โ is also known to exist at more moderate baryon densities. In particular, pairing of protons or neutrons inside large nuclei or neutron stars leads to superconductivity or superfluidity. Understanding the mechanism of nucleon pairing from the microscopic QCD theory may be as hard as understanding the mechanism for high-temperature superconductivity directly from the Hubbard model. Instead it is much more useful to employ a systematic low-energy effective theory whose parameters can be determined from the underlying microscopic physics. In nuclear physics effective field theory has recently led to some progress in describing the forces between nucleons in terms of just a few low-energy parameters , while phenomenological models involve a much larger number of adjustable parameters. Also steps towards describing nuclear matter with effective field theories have already been taken . The goal of the present paper is to develop a similar effective theory describing the interactions between the charge carriers in an antiferromagnet through magnon exchange. Remarkably, some physical phenomena that are practically inaccessible to microscopic Hubbard-type models even by numerical simulation can be tackled analytically in the effective field theory framework.
An ambitious goal of the effective theory approach is to systematically investigate possible mechanisms for the preformation of electron or hole pairs as a potential step towards understanding high-temperature superconductivity. It is an experimental fact that antiferromagnetism is destroyed before one enters the superconducting phase. How can magnon exchange then possibly provide a mechanism relevant for Cooper pair preformation? The destruction of antiferromagnetism just means the absence of infinite-range antiferromagnetic order. Antiferromagnetic correlations, although only of finite range, exist even in the superconducting phase. The finite correlation length implies that the magnons have developed a massgap, but they may still exist as relevant low-energy degrees of freedom. In particular, in $`2+1`$ dimensions, as a consequence of the Hohenberg-Mermin-Wagner-Coleman theorem, magnons pick up a mass that is exponentially small in the inverse temperature . The generation of the massgap is a non-perturbative phenomenon that is well within the applicability range of the effective theory, although infinite-range antiferromagnetic order exists only at zero temperature. Similarly, an effective theory for magnons and electrons or holes remains valid in the superconducting phase as long as the magnons remain among the lightest degrees of freedom. Again, this is similar to QCD where pions are not exactly massless either โ in that case as a result of explicit chiral symmetry breaking due to non-zero quark masses. Although pions are hence only pseudo-Nambu-Goldstone bosons, chiral perturbation theory remains perfectly well applicable.
The low-energy effective theory for magnons and charge carriers to be developed here is the condensed matter analog of baryon chiral perturbation theory in strong interaction physics . The effective theory is based on a non-linear realization of the spontaneously broken symmetry . The terms in the low-energy effective Lagrangian are organized according to the number of derivatives they contain. The lowest energy physics is dominated by the terms with the smallest number of derivatives, while effects at higher energies are taken into account systematically through higher-derivative terms. A key ingredient in constructing the effective Lagrangian are symmetry considerations. At a given order of the low-energy expansion, i.e. for a given number of derivatives, all terms consistent with the symmetries must be included in the effective Lagrangian, with a low-energy parameter that determines the strength of the corresponding interaction. For cuprates the most important symmetries are the $`SU(2)_s`$ spin symmetry which is spontaneously broken down to $`U(1)_s`$ in the antiferromagnetic phase, as well as the $`U(1)_Q`$ fermion number symmetry whose breakdown signals superconductivity. Other relevant symmetries include translation by one lattice spacing which changes the sign of the staggered magnetization, 90 degrees rotations and reflections of the square crystal lattice, as well as time-reversal. In addition to these generic symmetries of high-$`T_c`$ materials, the Hubbard model possesses an $`SU(2)_Q`$ symmetry discussed by Yang and Zhang which is a non-Abelian extension of the charge symmetry $`U(1)_Q`$. This symmetry is not expected to be present in generic cuprate materials, but may still be a relevant approximate symmetry in specific samples.
In this paper we ignore phonons, assuming that they do not play an important role for high-temperature superconductivity. For example, in the Hubbard model a rigid lattice which does not have its own physical degrees of freedom is put by hand. Of course, in the actual high-$`T_c`$ materials a crystal lattice arises as a result of the spontaneous breakdown of translation and Galilean (or more precisely Poincarรฉ) invariance. The corresponding Nambu-Goldstone bosons are the phonons. The role of phonons and their possible interplay with magnons can also be investigated systematically in the framework of low-energy effective field theory.
We also consider the coupling of antiferromagnets to external electromagnetic fields which can be used to probe the dynamics of magnons and electrons or holes. As first noted by Frรถhlich and Studer, in non-relativistic condensed matter external electromagnetic fields $`\stackrel{}{E}(x)`$ and $`\stackrel{}{B}(x)`$ enter the dynamics in the form of non-Abelian vector potentials for the $`SU(2)_s`$ spin symmetry . We use this observation to couple both the microscopic Hubbard model and the effective theory to external $`\stackrel{}{E}(x)`$ and $`\stackrel{}{B}(x)`$ fields. As discussed in detail in , the electromagnetic couplings are the condensed matter analog of the weak interactions in particle physics. These are described by an $`SU(2)_LU(1)_Y`$ gauge theory, which turns part of QCDโs global chiral symmetry into a gauge symmetry. Remarkably, the electromagnetic couplings of non-relativistic condensed matter are described by a local $`SU(2)_sU(1)_Q`$ symmetry which is the condensed matter analog of the $`SU(2)_LU(1)_Y`$ symmetry in particle physics. Some correspondences between QCD and antiferromagnets are summarized in table 1. Connections between QCD and condensed matter physics have also been discussed in .
The rest of this paper is organized as follows. Section 2 contains a symmetry analysis of the Hubbard model as a concrete example for an underlying microscopic system. In section 3 the effective theory for magnons is reviewed and the non-linear realization of the $`SU(2)_s`$ spin symmetry is constructed. In section 4 the Hubbard model is coupled to a magnon background field. In this way the fields of the effective theory inherit their transformation properties under the various symmetries from the underlying microscopic degrees of freedom. In section 5 the low-energy effective theory for magnons and charge carriers is developed and the leading terms in a systematic low-energy expansion of the effective action are constructed. This section also contains an application of the effective theory to the dispersion relations of charge carriers. Section 6 treats the $`t`$-$`J`$ model and its effective theory as a special case of systems with holes as the only charge carriers. In section 7 the Hubbard model as well as its effective theory are coupled to external electromagnetic fields. Finally, section 8 contains our conclusions, while some technical details are discussed in two appendices.
## 2 Symmetries of the Hubbard Model
In order to have a concrete microscopic system for which we will then construct a low-energy effective theory, we consider the Hubbard model. The Hubbard model just serves as one representative of a large class of systems, including the actual high-$`T_c`$ materials. Here it is essential that the Hubbard model shares important symmetries, e.g. an $`SU(2)_s`$ spin symmetry and a $`U(1)_Q`$ fermion number symmetry with these materials. In the Hubbard model at half-filling the $`U(1)_Q`$ symmetry even extends to an $`SU(2)_Q`$ symmetry. The $`SU(2)_Q`$ symmetry is not exact in actual materials, but may still be approximately realized and will also be investigated in the framework of the effective theory.
### 2.1 Hamiltonian and Generic Continuous Symmetries
The Hubbard model is defined by the Hamiltonian
$`H`$ $`=`$ $`t{\displaystyle \underset{x,i}{}}(c_x^{}c_{x+\widehat{i}}+c_{x+\widehat{i}}^{}c_x+c_x^{}c_{x+\widehat{i}}+c_{x+\widehat{i}}^{}c_x)`$ (2.1)
$`+U{\displaystyle \underset{x}{}}c_x^{}c_xc_x^{}c_x\mu ^{}{\displaystyle \underset{x}{}}(c_x^{}c_x+c_x^{}c_x).`$
Here $`x`$ denotes the sites of a 2-dimensional square lattice and $`\widehat{i}`$ is a vector of length $`a`$ (where $`a`$ is the lattice spacing) pointing in the $`i`$-direction. Furthermore, $`t`$ is the nearest-neighbor hopping parameter, while $`U>0`$ is the strength of the screened on-site Coulomb repulsion, and $`\mu ^{}`$ is the chemical potential for fermion number. The fermion creation and annihilation operators obey the standard anticommutation relations
$$\{c_{xs}^{},c_{ys^{}}\}=\delta _{xy}\delta _{ss^{}},\{c_{xs},c_{ys^{}}\}=\{c_{xs}^{},c_{ys^{}}^{}\}=0.$$
(2.2)
We also introduce the $`SU(2)_s`$ Pauli spinor
$$c_x=\left(\begin{array}{c}c_x\\ c_x\end{array}\right)$$
(2.3)
in terms of which (up to an irrelevant constant) the Hamiltonian takes the manifestly $`SU(2)_s`$-invariant form
$$H=t\underset{x,i}{}(c_x^{}c_{x+\widehat{i}}+c_{x+\widehat{i}}^{}c_x)+\frac{U}{2}\underset{x}{}(c_x^{}c_x1)^2\mu \underset{x}{}(c_x^{}c_x1).$$
(2.4)
Here $`\mu =\mu ^{}\frac{1}{2}U`$ is the chemical potential for the fermion number relative to half-filling, i.e. $`\mu =0`$ implies an average density of one fermion per lattice site. The corresponding $`U(1)_Q`$ symmetry is generated by the charge operator
$$Q=\underset{x}{}Q_x=\underset{x}{}(c_x^{}c_x1).$$
(2.5)
Again, we count fermion number relative to half-filling. The $`SU(2)_s`$ symmetry is generated by the total spin
$$\stackrel{}{S}=\underset{x}{}\stackrel{}{S}_x=\underset{x}{}c_x^{}\frac{\stackrel{}{\sigma }}{2}c_x,$$
(2.6)
where $`\stackrel{}{\sigma }`$ are the Pauli matrices. It is easy to see that the above Hamiltonian conserves both fermion number and spin, i.e. $`[H,Q]=[H,\stackrel{}{S}]=0`$, and that $`[Q,\stackrel{}{S}]=0`$. The infinitesimal generators $`\stackrel{}{S}`$ of $`SU(2)_s`$ (which obey the standard commutation relations $`[S_a,S_b]=i\epsilon _{abc}S_c`$) can be used to construct a unitary operator
$$V=\mathrm{exp}(i\stackrel{}{\eta }\stackrel{}{S}),$$
(2.7)
which implements the corresponding symmetry transformations in the Hilbert space of the theory. In particular, the transformed annihilation operators take the form
$$c_x^{}=V^{}c_xV=\mathrm{exp}(i\stackrel{}{\eta }\frac{\stackrel{}{\sigma }}{2})c_x=gc_x,g=\mathrm{exp}(i\stackrel{}{\eta }\frac{\stackrel{}{\sigma }}{2})SU(2)_s.$$
(2.8)
Similarly, the $`U(1)_Q`$ transformations are implemented by a unitary operator
$$W=\mathrm{exp}(i\omega Q),$$
(2.9)
such that
$${}_{}{}^{Q}c_{x}^{}=W^{}c_xW=\mathrm{exp}(i\omega )c_x,\mathrm{exp}(i\omega )U(1)_Q.$$
(2.10)
For large positive $`U`$, at half-filling, the repulsive Hubbard model reduces to the antiferromagnetic spin $`\frac{1}{2}`$ quantum Heisenberg model with the Hamiltonian
$$H=J\underset{x,i}{}\stackrel{}{S}_x\stackrel{}{S}_{x+\widehat{i}},$$
(2.11)
where the exchange coupling is given by $`J=2t^2/U`$. This follows to second order of perturbation theory in $`t/U`$. To leading order, i.e. completely ignoring the kinetic term proportional to $`t`$, there is an enormous number of degenerate ground states. Irrespective of spin, any state with exactly one fermion occupying each lattice site avoids the on-site Coulomb repulsion and thus represents a ground state for $`t=0`$. There is no correction at order $`t/U`$. In second order of degenerate perturbation theory, a spin can virtually hop to a neighboring site occupied by a fermion with opposite spin and then hop back. On the other hand, virtual hops to sites occupied by a fermion with the same spin orientation are forbidden by the Pauli principle. This favors antiparallel spins and leads to the antiferromagnetic Heisenberg model of eq.(2.11).
### 2.2 Discrete Symmetries
Since the Hubbard model at half-filling leads to antiferromagnetism, another important symmetry is translation by one lattice spacing (in the $`i`$-direction), which flips the sign of the staggered magnetization vector
$$\stackrel{}{M}_s=\underset{x}{}(1)^x\stackrel{}{S}_x.$$
(2.12)
The factor $`(1)^x=(1)^{(x_1+x_2)/a}`$ distinguishes between the sites of the even and odd sublattice. The points on the even sublattice $`A`$ have $`(1)^x=1`$ while the points on the odd sublattice $`B`$ have $`(1)^x=1`$. The displacement symmetry is generated by a unitary operator $`D`$ which acts as
$${}_{}{}^{D}c_{x}^{}=D^{}c_xD=c_{x+\widehat{i}},$$
(2.13)
and for which $`[H,D]=0`$. Obviously, both the $`U(1)_Q`$ and the $`SU(2)_s`$ symmetry commute with the displacement, i.e. $`[Q,D]=[\stackrel{}{S},D]=0`$. In the effective theory it will be useful to also consider a related symmetry $`D^{}`$ which combines $`D`$ with the spin rotation $`g=i\sigma _2`$. This symmetry acts as
$${}_{}{}^{D^{}}c_{x}^{}=D^{}c_xD^{}=(i\sigma _2)^Dc_x=(i\sigma _2)c_{x+\widehat{i}}.$$
(2.14)
Also note that $`[H,D^{}]=[D,D^{}]=[Q,D^{}]=0`$, but $`[\stackrel{}{S},D^{}]0`$.
In non-relativistic physics orbital angular momentum and spin are separately conserved and spin plays the role of an internal quantum number. Indeed, in the Hubbard model the $`SU(2)_s`$ spin symmetry is completely independent of the 90 degrees rotation invariance of the spatial lattice. The 90 degrees rotation $`O`$ acts on a spatial point $`x=(x_1,x_2)`$ as $`Ox=(x_2,x_1)`$. Under the symmetry $`O`$ the fermion operators transform as
$${}_{}{}^{O}c_{x}^{}=O^{}c_xO=c_{Ox}.$$
(2.15)
Parity turns $`x`$ into $`(x_1,x_2)`$ and is equivalent to a 180 degrees rotation in two dimensions. Hence, it is more useful to consider the spatial reflection $`R`$ at the $`x_1`$-axis which turns $`x`$ into $`Rx=(x_1,x_2)`$. Under this transformation the fermion operators transform as
$${}_{}{}^{R}c_{x}^{}=R^{}c_xR=c_{Rx}.$$
(2.16)
The reflection at the orthogonal $`x_2`$-axis is a combination of the reflection $`R`$ and the rotation $`O`$. One can also consider the reflection at an axis half between lattice points. This transformation is a combination of $`R`$ with the displacement symmetry $`D`$. Similarly, a reflection at a lattice diagonal is a combination of $`R`$ and $`O`$. Another important symmetry is time-reversal which is implemented by an antiunitary operator $`T`$.
It should be pointed out that, unlike the actual high-$`T_c`$ materials, the Hubbard model is not Galilean invariant: in the actual materials translation as well as Galilean invariance are spontaneously broken by the formation of the crystal lattice. The corresponding Nambu-Goldstone bosons are the phonons which are known to play a central role in ordinary low-$`T_c`$ superconductivity. In the Hubbard model, on the other hand, the lattice is imposed by hand, and thus translation and Galilean invariance are explicitly broken. In particular, phonons cannot arise because the lattice does not have its own physical degrees of freedom.
### 2.3 $`SU(2)_Q`$ Symmetry
As first noted by Yang and Zhang , at half-filling (i.e. for $`\mu =0`$) the Hubbard model possesses a non-Abelian extension $`SU(2)_Q`$ of the fermion number symmetry $`U(1)_Q`$ generated by
$$Q^+=\underset{x}{}(1)^xc_x^{}c_x^{},Q^{}=\underset{x}{}(1)^xc_xc_x,Q^3=\underset{x}{}\frac{1}{2}(c_x^{}c_x+c_x^{}c_x1)=\frac{1}{2}Q.$$
(2.17)
Writing $`Q^\pm =Q^1\pm iQ^2`$, it is straightforward to show that, for $`\mu =0`$, indeed $`[H,\stackrel{}{Q}]=0`$. Also the $`SU(2)_Q`$ symmetry commutes with the $`SU(2)_s`$ symmetry, i.e. $`[Q^a,S^b]=0`$, but it does not commute with the displacement symmetry because $`D^{}Q^\pm D=Q^\pm `$. For the same reason $`[\stackrel{}{Q},D^{}]0`$.
Introducing the $`SU(2)_Q`$ spinor
$$d_x=\left(\begin{array}{c}c_x\\ (1)^xc_x^{}\end{array}\right),$$
(2.18)
which obeys the standard anticommutation relations
$$\{d_{xa}^{},d_{yb}\}=\delta _{xy}\delta _{ab},\{d_{xa},d_{yb}\}=\{d_{xa}^{},d_{yb}^{}\}=0,$$
(2.19)
one writes
$$\stackrel{}{Q}=\underset{x}{}\stackrel{}{Q}_x=\underset{x}{}d_x^{}\frac{\stackrel{}{\sigma }}{2}d_x,$$
(2.20)
where $`\stackrel{}{\sigma }`$ are again the Pauli matrices now operating in $`SU(2)_Q`$ space. The infinitesimal generators $`\stackrel{}{Q}`$ of $`SU(2)_Q`$ can be used to construct a unitary operator
$$W=\mathrm{exp}(i\stackrel{}{\omega }\stackrel{}{Q}),$$
(2.21)
which implements the corresponding symmetry transformations in the Hilbert space of the theory. The transformed $`SU(2)_Q`$ spinors are then given by
$${}_{}{}^{\stackrel{}{Q}}d_{x}^{}=W^{}d_xW=\mathrm{exp}(i\stackrel{}{\omega }\frac{\stackrel{}{\sigma }}{2})d_x=\mathrm{\Omega }d_x,\mathrm{\Omega }=\mathrm{exp}(i\stackrel{}{\omega }\frac{\stackrel{}{\sigma }}{2})SU(2)_Q.$$
(2.22)
In terms of the $`SU(2)_Q`$ spinors the Hamiltonian takes the form
$$H=t\underset{x,i}{}(d_x^{}d_{x+\widehat{i}}+d_{x+\widehat{i}}^{}d_x)\frac{U}{2}\underset{x}{}(d_x^{}d_x1)^2\mu \underset{x}{}d_x^{}\sigma _3d_x.$$
(2.23)
The first two terms on the right-hand side are manifestly $`SU(2)_Q`$-invariant, while away from half-filling (i.e. for $`\mu 0`$) the chemical potential term explicitly breaks the $`SU(2)_Q`$ symmetry down to $`U(1)_Q`$.
Finally, we introduce a matrix-valued fermion operator
$$C_x=\left(\begin{array}{cc}c_x& (1)^xc_x^{}\\ c_x& (1)^xc_x^{}\end{array}\right),$$
(2.24)
which displays both the $`SU(2)_s`$ and the $`SU(2)_Q`$ symmetries in a compact form. The first column of $`C_x`$ is the $`SU(2)_s`$ spinor $`c_x`$, while the second column is another $`SU(2)_s`$ spinor which transforms exactly like $`c_x`$. The first row of $`C_x`$ is the $`SU(2)_Q`$ spinor $`d_x^T`$, while the second row is another $`SU(2)_Q`$ spinor which transforms exactly like $`d_x^T`$. Under combined $`SU(2)_s`$ and $`SU(2)_Q`$ transformations $`C_x`$ transforms as
$${}_{}{}^{\stackrel{}{Q}}C_{x}^{}=gC_x\mathrm{\Omega }^T.$$
(2.25)
Since the $`SU(2)_s`$ symmetry acts on the left while the $`SU(2)_Q`$ symmetry acts on the right, it is now manifest that the two symmetry operations commute. Under the displacement symmetry one obtains
$${}_{}{}^{D}C_{x}^{}=C_{x+\widehat{i}}\sigma _3.$$
(2.26)
The appearance of $`\sigma _3`$ on the right is due to the factor $`(1)^x`$ and confirms that the displacement symmetry commutes with all $`SU(2)_s`$ transformations, but only with the Abelian $`U(1)_Q`$ (and not with all $`SU(2)_Q`$) transformations. Similarly, under the symmetry $`D^{}`$ one finds
$${}_{}{}^{D^{}}C_{x}^{}=(i\sigma _2)C_{x+\widehat{i}}\sigma _3.$$
(2.27)
The Hamiltonian can now be expressed in a manifestly $`SU(2)_s`$-, $`U(1)_Q`$-, $`D`$-, and $`D^{}`$-invariant form
$$H=\frac{t}{2}\underset{x,i}{}\text{Tr}[C_x^{}C_{x+\widehat{i}}+C_{x+\widehat{i}}^{}C_x]+\frac{U}{12}\underset{x}{}\text{Tr}[C_x^{}C_xC_x^{}C_x]\frac{\mu }{2}\underset{x}{}\text{Tr}[C_x^{}C_x\sigma _3].$$
(2.28)
The chemical potential term is only $`U(1)_Q`$ invariant, while the other two terms are manifestly $`SU(2)_Q`$-invariant.
## 3 Effective Theory for Magnons
Before doping, the high-$`T_c`$ materials are quantum antiferromagnets in which the $`SU(2)_s`$ spin symmetry is spontaneously broken down to $`U(1)_s`$. The low-energy physics of antiferromagnets is dominated by the corresponding Nambu-Goldstone bosons โ the magnons. Chiral perturbation theory, which was originally developed for the Nambu-Goldstone pions of QCD, is a systematic low-energy expansion that has also been applied to magnons . In this section we review the basic features of magnon chiral perturbation theory. As a necessary prerequisite for the coupling of magnons to charge carriers, we also construct the non-linear realization of a spontaneously broken $`SU(2)_s`$ symmetry, which then appears as a local $`U(1)_s`$ symmetry in the unbroken subgroup. This is analogous to baryon chiral perturbation theory in which the spontaneously broken $`SU(2)_LSU(2)_R`$ chiral symmetry of QCD is implemented on the nucleon fields as a local $`SU(2)_{L=R}`$ transformation in the unbroken isospin subgroup.
### 3.1 Continuous Symmetries of the Effective Action
The undoped precursors of high-temperature layered cuprate superconductors are quantum antiferromagnets. At half-filling, also the Hubbard model displays antiferromagnetism. In these systems, at least at zero temperature, the spin rotational symmetry $`G=SU(2)_s`$ is spontaneously broken down to the subgroup $`H=U(1)_s`$ by the formation of a staggered magnetization. The $`U(1)_Q`$ symmetry, on the other hand, remains unbroken until one reaches the superconducting phase. In the Hubbard model even the $`SU(2)_Q`$ symmetry remains unbroken at half-filling but is explicitly broken down to $`U(1)_Q`$ for $`\mu 0`$. As a consequence of Goldstoneโs theorem, there are two massless bosons โ the antiferromagnetic spin-waves or magnons, which are described by a unit-vector field
$$\stackrel{}{e}(x)=(e_1(x),e_2(x),e_3(x)),\stackrel{}{e}(x)^2=1$$
(3.1)
in the coset space $`G/H=SU(2)_s/U(1)_s=S^2`$. Here $`x=(x_1,x_2,t)`$ denotes a point in Euclidean space-time. The vector $`\stackrel{}{e}(x)`$ describes the direction of the local staggered magnetization. The leading order terms in the Euclidean action of the low-energy effective theory for the magnons take the form
$$S[\stackrel{}{e}]=d^2x๐t\frac{\rho _s}{2}\left(_i\stackrel{}{e}_i\stackrel{}{e}+\frac{1}{c^2}_t\stackrel{}{e}_t\stackrel{}{e}\right).$$
(3.2)
The index $`i\{1,2\}`$ labels the two spatial directions, while the index $`t`$ refers to the Euclidean time-direction. The parameter $`\rho _s`$ is the spin stiffness and $`c`$ is the spin-wave velocity. For the antiferromagnetic Heisenberg model of eq.(2.11) these low-energy parameters have been determined very precisely in Monte Carlo calculations resulting in $`\rho _s=0.186(4)J`$, $`c=1.68(1)Ja`$, where $`J`$ is the exchange coupling of the Heisenberg model and $`a`$ is the lattice spacing. The leading terms in the magnon effective action are โPoincarรฉโ-invariant with the spin-wave velocity $`c`$ playing the role of the velocity of light. Consequently, antiferromagnetic magnons have a โrelativisticโ spectrum. The โPoincarรฉโ symmetry emerges only at low energies as a consequence of the discrete lattice rotation invariance. However, higher-derivative terms relevant at higher energies are in general not invariant.
In the following we prefer to work with an alternative representation of the magnon field using $`2\times 2`$ Hermitean projection matrices $`P(x)`$ that obey
$$P(x)^{}=P(x),\text{Tr}P(x)=1,P(x)^2=P(x),$$
(3.3)
and are given by
$$P(x)=\frac{1}{2}(\mathrm{๐ฃ}\mathrm{๐ฃ}+\stackrel{}{e}(x)\stackrel{}{\sigma })=\frac{1}{2}\left(\begin{array}{cc}1+e_3(x)& e_1(x)ie_2(x)\\ e_1(x)+ie_2(x)& 1e_3(x)\end{array}\right).$$
(3.4)
In the above $`CP(1)`$ language, the lowest-order effective action of eq.(3.2) takes the form
$$S[P]=d^2x๐t\rho _s\text{Tr}\left[_iP_iP+\frac{1}{c^2}_tP_tP\right].$$
(3.5)
This action is invariant under the global transformations $`gSU(2)_s`$ of eq.(2.8),
$$P(x)^{}=gP(x)g^{}.$$
(3.6)
Note that the magnon field $`P(x)`$ is invariant under the charge symmetries $`U(1)_Q`$ and $`SU(2)_Q`$, i.e. $`{}_{}{}^{\stackrel{}{Q}}P(x)=P(x)`$.
### 3.2 Discrete Symmetries of Magnon Fields
Under the displacement $`D`$ by one lattice spacing the staggered magnetization changes sign, i.e.
$$^D\stackrel{}{e}(x)=\stackrel{}{e}(x)^DP(x)=\mathrm{๐ฃ}\mathrm{๐ฃ}P(x).$$
(3.7)
Let us again combine $`D`$ with the spin rotation $`g=i\sigma _2`$, which results in the transformation $`D^{}`$ with
$${}_{}{}^{D^{}}P(x)=(i\sigma _2)^DP(x)(i\sigma _2)^{}=(i\sigma _2)[\mathrm{๐ฃ}\mathrm{๐ฃ}P(x)](i\sigma _2)^{}=P(x)^{},$$
(3.8)
reminiscent of charge conjugation in particle physics.
The Hubbard model is invariant under translations by an integer multiple of the lattice spacing. As we have seen, due to the antiferromagnetic order, the displacement $`D`$ by one lattice spacing (which connects the two sublattices $`A`$ and $`B`$) plays a special role. In particular, in the effective theory it manifests itself as an internal symmetry that changes the sign of $`\stackrel{}{e}(x)`$. Translations by an even number of lattice spacings (which do not mix the sublattices), on the other hand, manifest themselves as ordinary translations in the effective theory. It should be noted that in the effective theory one need not distinguish between the displacement symmetries $`D`$ for the two spatial directions, since they are related by an ordinary translation by two lattice spacings (one in the $`1`$\- and one in the $`2`$-direction).
When we decompose a space-time vector $`x=(x_1,x_2,t)`$ into its spatial and temporal components, the 90 degrees rotation $`O`$ acts on $`x`$ as $`Ox=(x_2,x_1,t)`$. Under the symmetry $`O`$ the magnon field transforms as
$${}_{}{}^{O}P(x)=P(Ox).$$
(3.9)
Similarly, under the spatial reflection $`R`$ at the $`x_1`$-axis, which turns $`x`$ into $`Rx=(x_1,x_2,t)`$, the magnon field transforms as
$${}_{}{}^{R}P(x)=P(Rx).$$
(3.10)
Had we not treated spin as an internal quantum number, it would also be directly affected by the spatial reflection. Since spin is a form of angular momentum, it transforms like the orbital angular momentum $`\stackrel{}{L}=\stackrel{}{r}\times \stackrel{}{p}`$ of a particle, which is a pseudo-vector and thus changes into $`{}_{}{}^{R}\stackrel{}{L}=(L_1,L_2,L_3)`$ under the reflection $`R`$. This is equivalent to a 180 degrees $`SU(2)_s`$ rotation around the $`2`$-direction. Since we treat $`SU(2)_s`$ as an exact internal symmetry, the pure spatial inversion $`R`$ (without 180 degrees rotation of the spin) is also a symmetry.
Another important symmetry is time-reversal $`T`$ which turns $`x=(x_1,x_2,t)`$ into $`Tx=(x_1,x_2,t)`$. In a Hamiltonian description time-reversal is represented by an antiunitary operator. Here we discuss time-reversal in the framework of the Euclidean path integral. Again, the spin transforms like the orbital angular momentum $`\stackrel{}{L}`$ of a particle. The momentum $`\stackrel{}{p}`$ changes sign under time-reversal and so does $`\stackrel{}{L}`$, i.e. $`{}_{}{}^{T}\stackrel{}{L}=\stackrel{}{L}`$.<sup>1</sup><sup>1</sup>1Note that $`\stackrel{}{L}`$ does not obey the angular momentum commutation relations. This is a consequence of the antiunitary nature of $`T`$ which does not represent an ordinary symmetry (implemented by a unitary transformation) in Hilbert space. Consequently, under $`T`$ the staggered magnetization vector (which is built from microscopic spins) transforms as
$${}_{}{}^{T}\stackrel{}{e}(x)=\stackrel{}{e}(Tx)^TP(x)=\mathrm{๐ฃ}\mathrm{๐ฃ}P(Tx)=^DP(Tx).$$
(3.11)
Hence, time-reversal is closely related to the displacement symmetry of eq.(3.7). Just like the displacement symmetry $`D`$, time-reversal is spontaneously broken in an antiferromagnet. However, in contrast to a ferromagnet, the combination $`TD`$ of time-reversal and the displacement symmetry remains unbroken. Previously we have combined the displacement symmetry $`D`$ with the $`SU(2)_s`$ spin rotation $`i\sigma _2`$ in order to obtain the unbroken symmetry $`D^{}`$. In order to obtain an unbroken variant $`T^{}`$ of time-reversal we now combine $`T`$ with the spin rotation $`i\sigma _2`$ which yields
$${}_{}{}^{T^{}}P(x)=(i\sigma _2)^TP(x)(i\sigma _2)^{}=(i\sigma _2)^DP(Tx)(i\sigma _2)^{}=^D^{}P(Tx).$$
(3.12)
### 3.3 Non-Linear Realization of the $`SU(2)_s`$ Symmetry
In order to couple electron or hole fields to the magnons one must construct a non-linear realization of the spontaneously broken $`SU(2)_s`$ symmetry which then manifests itself as a local symmetry in the unbroken $`U(1)_s`$ subgroup of $`SU(2)_s`$. This local transformation is constructed from the global transformation $`gSU(2)_s`$ as well as from the local magnon field $`P(x)`$ as follows: one first diagonalizes the magnon field by a unitary transformation $`u(x)SU(2)_s`$, i.e.
$$u(x)P(x)u(x)^{}=\frac{1}{2}(\mathrm{๐ฃ}\mathrm{๐ฃ}+\sigma _3)=\left(\begin{array}{cc}1& 0\\ 0& 0\end{array}\right),u_{11}(x)0.$$
(3.13)
Note that, due to its projector properties, $`P(x)`$ has eigenvalues 0 and 1. In order to make $`u(x)`$ uniquely defined, we demand that the element $`u_{11}(x)`$ is real and non-negative. Otherwise the diagonalizing matrix $`u(x)`$ would be defined only up to a $`U(1)_s`$ phase. Using eq.(3.4) and spherical coordinates for $`\stackrel{}{e}(x)`$, i.e.
$$\stackrel{}{e}(x)=(\mathrm{sin}\theta (x)\mathrm{cos}\phi (x),\mathrm{sin}\theta (x)\mathrm{sin}\phi (x),\mathrm{cos}\theta (x)),$$
(3.14)
one obtains
$`u(x)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2(1+e_3(x))}}}\left(\begin{array}{cc}1+e_3(x)& e_1(x)ie_2(x)\\ e_1(x)ie_2(x)& 1+e_3(x)\end{array}\right)`$ (3.17)
$`=`$ $`\left(\begin{array}{cc}\mathrm{cos}(\frac{1}{2}\theta (x))& \mathrm{sin}(\frac{1}{2}\theta (x))\mathrm{exp}(i\phi (x))\\ \mathrm{sin}(\frac{1}{2}\theta (x))\mathrm{exp}(i\phi (x))& \mathrm{cos}(\frac{1}{2}\theta (x))\end{array}\right).`$ (3.20)
Under a global $`SU(2)_s`$ transformation $`g`$ the diagonalizing field $`u(x)`$ transforms as
$$u(x)^{}=h(x)u(x)g^{},u_{11}(x)^{}0,$$
(3.21)
which implicitly defines the non-linear symmetry transformation
$$h(x)=\mathrm{exp}(i\alpha (x)\sigma _3)=\left(\begin{array}{cc}\mathrm{exp}(i\alpha (x))& 0\\ 0& \mathrm{exp}(i\alpha (x))\end{array}\right)U(1)_s.$$
(3.22)
The transformation $`h(x)`$ is uniquely defined since we demand that $`u_{11}(x)^{}`$ is again real and non-negative. Note that with this definition of $`h(x)`$ indeed
$$u(x)^{}P(x)^{}u(x)^{}=\frac{1}{2}(\mathrm{๐ฃ}\mathrm{๐ฃ}+\sigma _3).$$
(3.23)
Interestingly, the global $`SU(2)_s`$ transformation $`g`$ manifests itself in the form of a local transformation $`h(x)U(1)_s`$ which inherits its $`x`$-dependence from the magnon field $`P(x)`$.
We still need to show that the $`SU(2)_s`$ group structure $`g=g_2g_1`$ is inherited by the non-linear $`U(1)_s`$ realization, i.e. $`h(x)=h_2(x)h_1(x)`$. First, we perform the global $`SU(2)_s`$ transformation $`g_1`$, i.e.
$$P(x)^{}=g_1P(x)g_1^{},u(x)^{}=h_1(x)u(x)g_1^{},$$
(3.24)
which defines the non-linear realization $`h_1(x)`$. Then we perform the subsequent global transformation $`g_2`$ which defines the non-linear realization $`h_2(x)`$, i.e.
$`P(x)^{\prime \prime }=g_2P(x)^{}g_2^{}=g_2g_1P(x)(g_2g_1)^{}=gP(x)g^{},`$
$`u(x)^{\prime \prime }=h_2(x)u(x)^{}g_2^{}=h_2(x)h_1(x)u(x)(g_2g_1)^{}=h(x)u(x)g^{}.`$ (3.25)
This indeed implies the correct group structure $`h(x)=h_2(x)h_1(x)`$.
Under the displacement symmetry $`D`$ the sign-change of the staggered magnetization $`\stackrel{}{e}(x)`$ implies
$`{}_{}{}^{D}u(x)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2(1e_3(x))}}}\left(\begin{array}{cc}1e_3(x)& e_1(x)+ie_2(x)\\ e_1(x)+ie_2(x)& 1e_3(x)\end{array}\right)`$ (3.28)
$`=`$ $`\left(\begin{array}{cc}\mathrm{sin}(\frac{1}{2}\theta (x))& \mathrm{cos}(\frac{1}{2}\theta (x))\mathrm{exp}(i\phi (x))\\ \mathrm{cos}(\frac{1}{2}\theta (x))\mathrm{exp}(i\phi (x))& \mathrm{sin}(\frac{1}{2}\theta (x))\end{array}\right)`$ (3.31)
$`=`$ $`\tau (x)u(x),`$ (3.32)
where
$`\tau (x)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{e_1(x)^2+e_2(x)^2}}}\left(\begin{array}{cc}0& e_1(x)+ie_2(x)\\ e_1(x)+ie_2(x)& 0\end{array}\right)`$ (3.35)
$`=`$ $`\left(\begin{array}{cc}0& \mathrm{exp}(i\phi (x))\\ \mathrm{exp}(i\phi (x))& 0\end{array}\right).`$ (3.38)
Note that $`{}_{}{}^{D}\tau (x)=\tau (x)=\tau (x)^{}`$, such that
$${}_{}{}^{DD}u(x)=^D\tau (x)^Du(x)=\tau (x)^{}\tau (x)u(x)=u(x),$$
(3.39)
as one would expect for the displacement symmetry. It should also be noted that โ like the $`SU(2)_s`$ symmetry โ the displacement symmetry is also spontaneously broken and hence realized in a non-linear (i.e. magnon-field-dependent) manner. Similarly, under the displacement symmetry $`D^{}`$ one finds $`{}_{}{}^{D^{}}u(x)=h(x)^Du(x)g^{}`$ with $`g=i\sigma _2`$. For this particular $`g`$ the local transformation takes the form $`h(x)=(i\sigma _2)\tau (x)^{}`$, such that
$${}_{}{}^{D^{}}u(x)=u(x)^{}.$$
(3.40)
In contrast to the displacement symmetry $`D`$, the symmetry $`D^{}`$ is not spontaneously broken and is thus realized in a linear (i.e. magnon-field-independent) manner.
In the next step we consider the anti-Hermitean field
$$v_\mu (x)=u(x)_\mu u(x)^{},$$
(3.41)
which transforms under $`SU(2)_s`$ as
$$v_\mu (x)^{}=h(x)u(x)g^{}_\mu [gu(x)^{}h(x)^{}]=h(x)[v_\mu (x)+_\mu ]h(x)^{}.$$
(3.42)
Writing
$$v_\mu (x)=iv_\mu ^a(x)\sigma _a=i\left(\begin{array}{cc}v_\mu ^3(x)& v_\mu ^+(x)\\ v_\mu ^{}(x)& v_\mu ^3(x)\end{array}\right),v_\mu ^\pm (x)=v_\mu ^1(x)iv_\mu ^2(x)$$
(3.43)
and using eq.(3.22) this implies
$$v_\mu ^3(x)^{}=v_\mu ^3(x)_\mu \alpha (x),v_\mu ^\pm (x)^{}=\mathrm{exp}(\pm 2i\alpha (x))v_\mu ^\pm (x).$$
(3.44)
Hence, $`v_\mu ^3(x)`$ transforms like an Abelian gauge field for $`U(1)_s`$, while $`v_\mu ^\pm (x)`$ represent vector fields โchargedโ under $`U(1)_s`$. For later convenience we also introduce the Hermitean charged vector field
$$V_\mu (x)=v_\mu ^1(x)\sigma _1+v_\mu ^2(x)\sigma _2=v_\mu ^+(x)\sigma _++v_\mu ^{}(x)\sigma _{}=\left(\begin{array}{cc}0& v_\mu ^+(x)\\ v_\mu ^{}(x)& 0\end{array}\right),$$
(3.45)
where $`\sigma _\pm =\frac{1}{2}(\sigma _1\pm i\sigma _2)`$ are raising and lowering operators of spin. Under the $`SU(2)_s`$ symmetry the charged vector field transforms as
$$V_\mu (x)^{}=h(x)V_\mu (x)h(x)^{}.$$
(3.46)
The magnon action can also be written as
$$S[v_\mu ]=d^2x๐t2\rho _s\left(v_i^+v_i^{}+\frac{1}{c^2}v_t^+v_t^{}\right)=d^2x๐t\rho _s\text{Tr}\left[V_i^{}V_i+\frac{1}{c^2}V_t^{}V_t\right].$$
(3.47)
It should be pointed out that the fields $`v_\mu ^a(x)`$ do not represent independent degrees of freedom, but are composed of magnon fields. In particular, what looks like a mass term for a charged vector field is indeed just the kinetic term of a massless Nambu-Goldstone boson.
### 3.4 Discrete Symmetries of Composite Fields
Under the displacement symmetry $`D`$ the composite vector field transforms as
$`{}_{}{}^{D}v_{\mu }^{}(x)=\tau (x)[v_\mu (x)+_\mu ]\tau (x)^{}^Dv_\mu ^3(x)=v_\mu ^3(x)+_\mu \phi (x),`$
$`{}_{}{}^{D}v_{\mu }^{\pm }(x)=\mathrm{exp}(2i\phi (x))v_\mu ^{}(x),^DV_\mu (x)=\tau (x)V_\mu (x)\tau (x)^{}.`$ (3.48)
Similarly, under the symmetry $`D^{}`$ one finds
$`{}_{}{}^{D^{}}v_{\mu }^{}(x)=v_\mu (x)^{}^D^{}v_\mu ^3(x)=v_\mu ^3(x),`$
$`{}_{}{}^{D^{}}v_{\mu }^{\pm }(x)=v_\mu ^{}(x),^D^{}V_\mu (x)=V_\mu (x)^{}.`$ (3.49)
This is exactly how an ordinary non-Abelian gauge field behaves under charge conjugation in particle physics.
Under the 90 degrees spatial rotation $`O`$ the composite field $`v_\mu (x)`$ transforms as
$${}_{}{}^{O}v_{i}^{}(x)=\epsilon _{ij}v_j(Ox),^Ov_t(x)=v_t(Ox),$$
(3.50)
while under the reflection $`R`$ one obtains
$${}_{}{}^{R}v_{1}^{}(x)=v_1(Rx),^Rv_2(x)=v_2(Rx),^Rv_t(x)=v_t(Rx).$$
(3.51)
Finally, under the time-reversal symmetry $`T`$ the field $`v_\mu `$ transforms as
$`{}_{}{}^{T}v_{i}^{}(x)=^Dv_i(Tx),^Tv_t(x)=^Dv_t(Tx)`$
$`{}_{}{}^{T}v_{i}^{3}(x)=v_i^3(Tx)+_i\phi (Tx),^Tv_t^3(x)=v_t^3(Tx)_t\phi (Tx),`$
$`{}_{}{}^{T}v_{i}^{\pm }(x)=\mathrm{exp}(2i\phi (Tx))v_i^{}(Tx),^Tv_t^\pm (x)=\mathrm{exp}(2i\phi (Tx))v_t^{}(Tx),`$
$`{}_{}{}^{T}V_{i}^{}(x)=\tau (Tx)V_i(Tx)\tau (Tx)^{},^TV_t(x)=\tau (Tx)V_t(Tx)\tau (Tx)^{},`$ (3.52)
and under its unbroken variant $`T^{}`$ one finds
$`{}_{}{}^{T^{}}v_{i}^{}(x)=^D^{}v_i(Tx),^T^{}v_t(x)=^D^{}v_t(Tx),`$
$`{}_{}{}^{T^{}}v_{i}^{3}(x)=v_i^3(Tx),^T^{}v_t^3(x)=v_t^3(Tx),`$
$`{}_{}{}^{T^{}}v_{i}^{\pm }(x)=v_i^{}(Tx),^T^{}v_t^\pm (x)=v_t^{}(Tx),`$
$`{}_{}{}^{T^{}}V_{i}^{}(x)=V_i(Tx)^T,^T^{}V_t(x)=V_t(Tx)^T.`$ (3.53)
Note that the upper index $`T`$ on the right denotes transpose, while on the left it denotes time-reversal. The above relations are equivalent to time-reversal of an ordinary non-Abelian gauge field.
### 3.5 Alternative Representation of Magnon Fields
We have used two equivalent representations of the magnon field in terms of the unit-vector $`\stackrel{}{e}(x)`$ and in terms of the projection matrix $`P(x)`$. There is a third equivalent representation in terms of a complex doublet
$`z(x)=\left(\begin{array}{c}z_1(x)\\ z_2(x)\end{array}\right),z(x)^{}=(z_1(x)^{},z_2(x)^{}),`$ (3.56)
$`z(x)^{}z(x)=|z_1(x)|^2+|z_2(x)|^2=1,`$ (3.57)
which is related to the other two representations by
$`\stackrel{}{e}(x)=z(x)^{}\stackrel{}{\sigma }z(x)`$
$`e_1(x)=z_1(x)^{}z_2(x)+z_2(x)^{}z_1(x),`$
$`e_2(x)=i[z_2(x)^{}z_1(x)z_1(x)^{}z_2(x)],`$
$`e_3(x)=|z_1(x)|^2|z_2(x)|^2,`$
$`P(x)=z(x)z(x)^{}=\left(\begin{array}{cc}|z_1(x)|^2& z_1(x)z_2(x)^{}\\ z_2(x)z_1(x)^{}& |z_2(x)|^2\end{array}\right).`$ (3.60)
The field $`z(x)`$ is defined in terms of $`\stackrel{}{e}(x)`$ or $`P(x)`$ only up to a $`U(1)_s`$ gauge transformation
$$z(x)^{}=\mathrm{exp}(i\beta (x))z(x).$$
(3.61)
It is therefore necessary to also introduce the auxiliary real-valued $`U(1)_s`$ gauge field
$$a_\mu (x)=\frac{1}{2i}[z(x)^{}_\mu z(x)_\mu z(x)^{}z(x)],$$
(3.62)
which under the symmetry of eq.(3.61) transforms as
$$a_\mu (x)^{}=a_\mu (x)+_\mu \beta (x).$$
(3.63)
The complex doublet $`z(x)`$ is closely related to the field $`u(x)`$. Fixing the gauge freedom of eq.(3.61) such that $`z_1(x)`$ is real and non-negative, it is easy to show that
$$u(x)=\left(\begin{array}{cc}z_1(x)& z_2(x)^{}\\ z_2(x)& z_1(x)\end{array}\right),v_\mu ^3(x)=a_\mu (x).$$
(3.64)
Hence, the description in terms of complex doublets $`z(x)`$ and an additional auxiliary gauge field $`a_\mu (x)`$ is physically equivalent to what we described before. It should again be pointed out that $`a_\mu (x)`$ (or equivalently $`v_\mu ^3(x)`$) does not represent a dynamical Abelian gauge field, but is simply a composite field constructed from the underlying magnon field $`P(x)`$.
### 3.6 Baby-Skyrmions
It is interesting to note that magnon fields support topological solitons known as baby-Skyrmions โ a lower-dimensional variant of the Skyrme soliton which represents a baryon in the low-energy pion effective theory for QCD . Baby-Skyrmions are solitons whose topological charge
$$B=\frac{1}{8\pi }d^2x\epsilon _{ij}\stackrel{}{e}(_i\stackrel{}{e}\times _j\stackrel{}{e}),$$
(3.65)
defined at every instant in time, is an element of the homotopy group $`\mathrm{\Pi }_2[S^2]=๐น๐น`$. The corresponding topological current
$$j_\mu (x)=\frac{1}{8\pi }\epsilon _{\mu \nu \rho }\stackrel{}{e}(x)[_\nu \stackrel{}{e}(x)\times _\rho \stackrel{}{e}(x)]$$
(3.66)
is conserved, i.e. $`_\mu j_\mu =0`$, independent of the equations of motion. Baby-Skyrmions are massive excitations inaccessible to the systematic low-energy expansion of chiral perturbation theory. Still, the existence of the conserved current $`j_\mu (x)`$ may have physical consequences even for the pure magnon dynamics.
Under the various symmetries the topological current transforms as
$`SU(2)_s:`$ $`j_\mu (x)^{}=j_\mu (x),`$
$`SU(2)_Q:`$ $`{}_{}{}^{\stackrel{}{Q}}j_{\mu }^{}(x)=j_\mu (x),`$
$`D:`$ $`{}_{}{}^{D}j_{\mu }^{}(x)=j_\mu (x),`$
$`D^{}:`$ $`{}_{}{}^{D^{}}j_{\mu }^{}(x)=j_\mu (x),`$
$`O:`$ $`{}_{}{}^{O}j_{t}^{}(x)=j_t(Ox),^Oj_i(x)=\epsilon _{ij}j_j(Ox),`$
$`R:`$ $`{}_{}{}^{R}j_{t}^{}(x)=j_t(Rx),^Rj_1(x)=j_1(Rx),^Rj_2(x)=j_2(Rx),`$
$`T:`$ $`{}_{}{}^{T}j_{t}^{}(x)=j_t(Tx),^Tj_1(x)=j_1(Tx),^Tj_2(x)=j_2(Rx),`$
$`T^{}:`$ $`{}_{}{}^{T^{}}j_{t}^{}(x)=j_t(Tx),^T^{}j_1(x)=j_1(Tx),^T^{}j_2(x)=j_2(Tx).`$ (3.67)
One might be tempted to add a term $`j_\mu (x)v_\mu ^3(x)`$ to the magnon Lagrangian because this is how an Abelian gauge field couples to a conserved current. Indeed, this term is invariant under $`SU(2)_s`$, $`SU(2)_Q`$, $`D`$, $`D^{}`$, and $`O`$. However, it violates the reflection and time-reversal symmetries $`R`$, $`T`$, and $`T^{}`$ and is hence forbidden.
There is another non-trivial homotopy group, $`\mathrm{\Pi }_3[S^2]=๐น๐น`$, which is relevant for baby-Skyrmions. It implies that space-time-dependent magnon fields fall into distinct topological classes characterized by the Hopf number $`H[\stackrel{}{e}]\mathrm{\Pi }_3[S^2]=๐น๐น`$. In $`2+1`$ dimensions baby-Skyrmions can be quantized as anyons characterized by a statistics angle $`\theta `$ . The cases $`\theta =0`$ and $`\theta =\pi `$ correspond to bosons and fermions, respectively. Including the Hopf term, the magnon path integral takes the form
$$Z=๐\stackrel{}{e}\mathrm{exp}(S[\stackrel{}{e}])\mathrm{exp}(i\theta H[\stackrel{}{e}]).$$
(3.68)
The Hopf term also changes sign under $`R`$, $`T`$, and $`T^{}`$. Hence, $`\mathrm{exp}(i\theta H[\stackrel{}{e}])`$ is invariant only if $`\theta `$ is 0 or $`\pi `$. Consequently, in an antiferromagnet with exact $`R`$, $`T`$, or $`T^{}`$ symmetries baby-Skyrmions can only be quantized as bosons or fermions. For the antiferromagnetic quantum Heisenberg model it has been argued that no Hopf term is generated . Hence, in that case the baby-Skyrmions should be bosons.
## 4 The Hubbard Model in a Magnon Background Field
The half-filled ground state of the Hubbard model plays a similar role as the Dirac sea in a relativistic quantum field theory. In particular, any fermion added to a half-filled state will be denoted as an electron, while any fermion removed from such a state represents a hole. In this section we couple a background magnon field to the microscopic degrees of freedom of the Hubbard model. In this way composite operators are constructed which transform exactly like the fields of the effective theory. Hence, the effective fields inherit their transformation properties under symmetry operations from the Hubbard model degrees of freedom.
### 4.1 Fermion Operators in a Magnon Background Field
In order to analyze the transformation properties of the electron and hole fields, as an intermediate step between the microscopic and effective descriptions, we first add a continuum magnon background field $`P(x)`$ to the Hubbard model by hand. The corresponding diagonalizing unitary matrix field $`u(x)`$ is used to turn the matrix-valued Hubbard model operator $`C_x`$ of eq.(2.24) into new operators $`\mathrm{\Psi }_x^A`$ and $`\mathrm{\Psi }_x^B`$ defined on the even and odd sublattices, respectively
$`\mathrm{\Psi }_x^A=u(x)C_x=u(x)\left(\begin{array}{cc}c_x& c_x^{}\\ c_x& c_x^{}\end{array}\right)=\left(\begin{array}{cc}\psi _{x+}^A& \psi _x^A\\ \psi _x^A& \psi _{x+}^A\end{array}\right),xA,`$ (4.5)
$`\mathrm{\Psi }_x^B=u(x)C_x=u(x)\left(\begin{array}{cc}c_x& c_x^{}\\ c_x& c_x^{}\end{array}\right)=\left(\begin{array}{cc}\psi _{x+}^B& \psi _x^B\\ \psi _x^B& \psi _{x+}^B\end{array}\right),xB.`$ (4.10)
In order to achieve a consistent representation of the underlying antiferromagnetic structure, it is unavoidable to explicitly split the degrees of freedom according to their location on sublattice $`A`$ or $`B`$. In this context it may be interesting to consider the electron-hole representation of the Hubbard model operators discussed in appendix A. The operators $`\psi _{x\pm }^{A,B}`$ obey standard anticommutation relations. It should be noted that here the continuum field $`u(x)`$ is evaluated only at discrete lattice points $`x`$.
The new lattice operators inherit their transformation properties from the operators of the Hubbard model. According to eqs.(3.21) and (2.8), under the $`SU(2)_s`$ symmetry one obtains
$$\mathrm{\Psi }_{x}^{A,B}{}_{}{}^{}=u(x)^{}C_x^{}=h(x)u(x)g^{}gC_x=h(x)\mathrm{\Psi }_x^{A,B}.$$
(4.11)
In components this relation takes the form
$$\psi _{x\pm }^{A,B}{}_{}{}^{}=\mathrm{exp}(\pm i\alpha (x))\psi _{x\pm }^{A,B}.$$
(4.12)
The components $`\psi _{x\pm }^{A,B}`$ do not simply correspond to spin up and spin down with respect to an arbitrarily chosen global quantization axis. Instead they correspond to spin parallel ($`+`$) or antiparallel ($``$) to the local staggered magnetization. This follows from considering global symmetry transformations $`gU(1)_s`$ in the unbroken subgroup of $`SU(2)_s`$ which describe rotations around the spontaneously selected direction of the staggered magnetization vector. In that case, according to eq.(3.21), $`h(x)=g`$ becomes a global transformation as well and eq.(4.12) shows that $`\psi _{x\pm }^{A,B}`$ indeed has spin parallel or antiparallel to the direction of the staggered magnetization.
Similarly, under the $`SU(2)_Q`$ symmetry one obtains
$${}_{}{}^{\stackrel{}{Q}}\mathrm{\Psi }_{x}^{A,B}=^\stackrel{}{Q}u(x)^\stackrel{}{Q}C_x=u(x)C_x\mathrm{\Omega }^T=\mathrm{\Psi }_x^{A,B}\mathrm{\Omega }^T.$$
(4.13)
In particular, under the $`U(1)_Q`$ subgroup of $`SU(2)_Q`$ the components transform as
$${}_{}{}^{Q}\psi _{x\pm }^{A,B}=\mathrm{exp}(i\omega )\psi _{x\pm }^{A,B}.$$
(4.14)
Under the displacement symmetry the new operators transform as
$${}_{}{}^{D}\mathrm{\Psi }_{x}^{A,B}=^Du(x+\widehat{i})C_{x+\widehat{i}}\sigma _3=\tau (x+\widehat{i})u(x+\widehat{i})C_{x+\widehat{i}}\sigma _3=\tau (x+\widehat{i})\mathrm{\Psi }_{x+\widehat{i}}^{B,A}\sigma _3,$$
(4.15)
where $`\tau (x)`$ is the field introduced in eq.(3.28). Expressed in components this implies
$${}_{}{}^{D}\psi _{x\pm }^{A,B}=\mathrm{exp}(i\phi (x+\widehat{i}))\psi _{x+\widehat{i},}^{B,A}.$$
(4.16)
Similarly, under the symmetry $`D^{}`$ one finds
$${}_{}{}^{D^{}}\mathrm{\Psi }_{x}^{A,B}=^D^{}u(x+\widehat{i})(i\sigma _2)C_{x+\widehat{i}}\sigma _3=u(x+\widehat{i})^{}(i\sigma _2)C_{x+\widehat{i}}\sigma _3=(i\sigma _2)\mathrm{\Psi }_{x+\widehat{i}}^{B,A}\sigma _3.$$
(4.17)
Here we have used $`u(x+\widehat{i})^{}(i\sigma _2)=(i\sigma _2)u(x+\widehat{i})`$. Again, expressed in components this relation takes the form
$${}_{}{}^{D^{}}\psi _{x\pm }^{A,B}=\pm \psi _{x+\widehat{i},}^{B,A}.$$
(4.18)
We have seen before that the symmetry $`D^{}`$ acts on the composite field $`v_\mu (x)`$ exactly like charge conjugation in particle physics. However, it should be noted that $`D^{}`$ acts on the electron and hole fields in a different way than the usual charge conjugation of a relativistic Dirac fermion which interchanges electrons and positrons. In particular, $`D^{}`$ does not interchange electrons and holes. Instead, it flips the spin of both electrons and holes from $`+`$ to $``$ and vice versa. Indeed, the spin is the โchargeโ that couples to the composite gauge field of eq.(3.41) constructed from the magnon field.
In the condensed matter literature on high-temperature superconductivity the concept of spin-charge separation (whose existence is established for some systems in one spatial dimension) has often been invoked. The idea is that there may be quasi-particles โ so-called holons โ which carry charge but no spin, as well as so-called spinons which are neutral and carry spin $`\frac{1}{2}`$. In order to avoid confusion between holons and the holes of our effective theory, we like to make a few comments: one might think that the fermion operator $`\mathrm{\Psi }_x^{A,B}`$ does not carry spin since it does not transform with the global spin transformation $`gSU(2)_s`$. However, the spin symmetry is non-linearly realized and hence the fermion operator transforms with the local $`h(x)U(1)_s`$. Consequently, $`\mathrm{\Psi }_x^{A,B}`$ still carries spin and hence does not represent a holon. It should also be pointed out that in the weakly coupled effective theory of magnons and holes there are no linearly confining forces that could form a spinless holon out of $`\mathrm{\Psi }_x^{A,B}`$ and the magnon field $`z(x)`$ of eq.(3.56).
### 4.2 Formal Continuum Limit of the Hubbard Model in a Magnon Background Field
In terms of the new operators the Hubbard model Hamiltonian takes the form
$`H`$ $`=`$ $`{\displaystyle \frac{t}{2}}{\displaystyle \underset{xA,i}{}}\text{Tr}[\mathrm{\Psi }_x^A๐ฑ_{x,i}\mathrm{\Psi }_{x+\widehat{i}}^B+\mathrm{\Psi }_{x+\widehat{i}}^B๐ฑ_{x,i}^{}\mathrm{\Psi }_x^A]`$ (4.19)
$`{\displaystyle \frac{t}{2}}{\displaystyle \underset{xB,i}{}}\text{Tr}[\mathrm{\Psi }_x^B๐ฑ_{x,i}\mathrm{\Psi }_{x+\widehat{i}}^A+\mathrm{\Psi }_{x+\widehat{i}}^A๐ฑ_{x,i}^{}\mathrm{\Psi }_x^B]`$
$`+{\displaystyle \frac{U}{12}}{\displaystyle \underset{xA}{}}\text{Tr}[\mathrm{\Psi }_x^A\mathrm{\Psi }_x^A\mathrm{\Psi }_x^A\mathrm{\Psi }_x^A]+{\displaystyle \frac{U}{12}}{\displaystyle \underset{xB}{}}\text{Tr}[\mathrm{\Psi }_x^B\mathrm{\Psi }_x^B\mathrm{\Psi }_x^B\mathrm{\Psi }_x^B]`$
$`{\displaystyle \frac{\mu }{2}}{\displaystyle \underset{xA}{}}\text{Tr}[\mathrm{\Psi }_x^A\mathrm{\Psi }_x^A\sigma _3]{\displaystyle \frac{\mu }{2}}{\displaystyle \underset{xB}{}}\text{Tr}[\mathrm{\Psi }_x^B\mathrm{\Psi }_x^B\sigma _3],`$
where we have introduced the parallel transporter
$$๐ฑ_{x,i}=u(x)u(x+\widehat{i})^{}SU(2)_s,$$
(4.20)
which transforms under $`SU(2)_s`$ as
$$๐ฑ_{x,i}^{}=h(x)๐ฑ_{x,i}h(x+\widehat{i})^{}.$$
(4.21)
For smooth magnon fields we can put
$`u(x)=u(x+{\displaystyle \frac{\widehat{i}}{2}}){\displaystyle \frac{a}{2}}_iu(x+{\displaystyle \frac{\widehat{i}}{2}})+{\displaystyle \frac{a^2}{8}}_i^2u(x+{\displaystyle \frac{\widehat{i}}{2}})+๐ช(a^3),`$
$`u(x+\widehat{i})=u(x+{\displaystyle \frac{\widehat{i}}{2}})+{\displaystyle \frac{a}{2}}_iu(x+{\displaystyle \frac{\widehat{i}}{2}})+{\displaystyle \frac{a^2}{8}}_i^2u(x+{\displaystyle \frac{\widehat{i}}{2}})+๐ช(a^3),`$ (4.22)
where $`a`$ is the lattice spacing. Similar expressions hold for the other fields. Using the unitarity of $`u(x+\frac{\widehat{i}}{2})`$ one can show that the lattice parallel transporter reduces to
$$๐ฑ_{x,i}=\mathrm{๐ฃ}\mathrm{๐ฃ}+av_i(x+\frac{\widehat{i}}{2})+\frac{a^2}{2}v_i(x+\frac{\widehat{i}}{2})^2+๐ช(a^3),$$
(4.23)
with $`v_i(x)`$ given by eq.(3.41). Note that both the continuum field $`v_i(x)`$ and the lattice parallel transporter field $`๐ฑ_{x,i}`$ transform locally only with the unbroken $`U(1)_s`$ subgroup and not with the full $`SU(2)_s`$ symmetry.
In the continuum limit we make the replacements
$$\underset{xA}{},\underset{xB}{}\frac{1}{2a^2}d^2x,\mathrm{\Psi }_x^{A,B}\sqrt{2}a\mathrm{\Psi }^{A,B}(x).$$
(4.24)
The factor $`\frac{1}{2}`$ in front of the integral accounts for the fact that each sublattice covers only half of the space. Similarly the factor $`\sqrt{2}a`$ in the definition of the continuum field $`\mathrm{\Psi }^{A,B}(x)`$ arises because there is only one degree of freedom of a given type $`A`$ or $`B`$ per area $`2a^2`$. The components $`\psi _\pm ^{A,B}(x)`$ of $`\mathrm{\Psi }^{A,B}(x)`$ again obey standard anticommutation relations, however, with the Dirac $`\delta `$-function of the continuum theory instead of the Kronecker $`\delta `$-function of the lattice. It should be noted that, due to the antiferromagnetic order, the number of degrees of freedom per continuum point is twice as large as the number per lattice point. Taking the formal continuum limit $`a0`$ (and ignoring an irrelevant constant) the Hamiltonian of eq.(4.19) takes the form
$`H`$ $`=`$ $`{\displaystyle }d^2x\{M\text{Tr}[\mathrm{\Psi }^A\mathrm{\Psi }^B]+{\displaystyle \frac{1}{2M^{}}}\text{Tr}[D_i\mathrm{\Psi }^AD_i\mathrm{\Psi }^B]`$
$`+iK\text{Tr}[D_i\mathrm{\Psi }^AV_i\mathrm{\Psi }^B+D_i\mathrm{\Psi }^BV_i\mathrm{\Psi }^A]+N\text{Tr}[\mathrm{\Psi }^AV_iV_i\mathrm{\Psi }^B]`$
$`+{\displaystyle \frac{G}{12}}\text{Tr}[\mathrm{\Psi }^A\mathrm{\Psi }^A\mathrm{\Psi }^A\mathrm{\Psi }^A+\mathrm{\Psi }^B\mathrm{\Psi }^B\mathrm{\Psi }^B\mathrm{\Psi }^B]{\displaystyle \frac{\mu }{2}}\text{Tr}[\mathrm{\Psi }^A\mathrm{\Psi }^A\sigma _3+\mathrm{\Psi }^B\mathrm{\Psi }^B\sigma _3]\}.`$
(4.25)
It should be noted that, due to the structure of $`\mathrm{\Psi }^{A,B}(x)`$, the individual terms are Hermitean. In the above expression $`V_\mu (x)`$ is the field defined in eq.(3.45) and the covariant derivatives are given by
$`D_\mu \mathrm{\Psi }^{A,B}(x)=(_\mu +iv_\mu ^3(x)\sigma _3)\mathrm{\Psi }^{A,B}(x),`$
$`D_\mu \mathrm{\Psi }^{A,B}(x)=[D_\mu \mathrm{\Psi }^{A,B}(x)]^{}=_\mu \mathrm{\Psi }^{A,B}(x)\mathrm{\Psi }^{A,B}(x)iv_\mu ^3(x)\sigma _3.`$ (4.26)
In terms of the fundamental parameters $`t`$ and $`U`$ and the lattice spacing $`a`$ of the Hubbard model one obtains
$$M=4t,M^{}=\frac{1}{2ta^2},K=ta^2,N=ta^2,G=2Ua^2.$$
(4.27)
It should be noted that (in contrast to a relativistic theory) the kinetic mass $`M^{}`$ is in general different from the rest mass $`M`$. The Hamiltonian from above resembles some (but not all) terms in the action of the effective theory to be constructed below. However, the coupling constants resulting from the formal continuum limit get renormalized and will hence be replaced by a priori unknown low-energy parameters in the effective action. The values of the low-energy parameters can be determined in experiments with cuprate materials or through numerical simulations of a microscopic Hubbard-type model.
## 5 Effective Theory for Magnons and Charge Carriers
The low-energy effective theory for magnons is analogous to chiral perturbation theory for pions in QCD. In QCD the baryon number $`B`$ is a conserved quantity. Thus one can investigate the low-energy QCD dynamics separately in each baryon number sector. Ordinary chiral perturbation theory operates in the $`B=0`$ sector. The low-energy physics in the $`B=1`$ sector involves a single nucleon interacting with soft pions. The low-energy effective theory describing these dynamics is known as baryon chiral perturbation theory . Similar effective theories have been constructed for the $`B=2`$ and $`B=3`$ sectors in the context of nuclear physics. Even nuclear matter (i.e. a system with non-zero baryon density) has been studied with effective theories . The condensed matter analog of baryon number is electron (or hole) number (or equivalently electric charge) which is obviously also conserved. In analogy to QCD it is hence possible to construct a low-energy effective theory describing the interactions of soft magnons with charge carriers doped into an antiferromagnet. Most high-$`T_c`$ materials result by hole-doping of quantum antiferromagnets, but the effective theory also applies to electron-doping. The key observation is that the spontaneously broken $`SU(2)_s`$ spin symmetry is non-linearly realized on the electron or hole fields and appears as a local $`U(1)_s`$ symmetry in the unbroken subgroup. This is analogous to baryon chiral perturbation theory in which the spontaneously broken $`SU(2)_LSU(2)_R`$ chiral symmetry of QCD is implemented on the nucleon fields as a local $`SU(2)_{L=R}`$ transformation in the unbroken isospin subgroup.
### 5.1 Effective Fields for Charge Carriers
In the low-energy effective theory we will use a Euclidean path integral description instead of the Hamiltonian description used in the Hubbard model. Consequently, the Hermitean conjugate lattice operators $`\psi _{x\pm }^{A,B}`$ are then replaced by Grassmann numbers $`\psi _\pm ^{A,B}(x)`$ which are completely independent of $`\psi _\pm ^{A,B}(x)`$. Therefore, in the effective theory the electron and hole fields are represented by eight independent Grassmann numbers $`\psi _\pm ^{A,B}(x)`$ and $`\psi _\pm ^{A,B}(x)`$ which can be combined to
$$\mathrm{\Psi }^A(x)=\left(\begin{array}{cc}\psi _+^A(x)& \psi _{}^A(x)\\ \psi _{}^A(x)& \psi _+^A(x)\end{array}\right),\mathrm{\Psi }^B(x)=\left(\begin{array}{cc}\psi _+^B(x)& \psi _{}^B(x)\\ \psi _{}^B(x)& \psi _+^B(x)\end{array}\right).$$
(5.1)
In order to avoid confusion with relativistic theories, we do not denote the conjugate fields by $`\overline{\psi }_\pm ^{A,B}(x)`$. For notational convenience we also introduce the fields
$$\mathrm{\Psi }^A(x)=\left(\begin{array}{cc}\psi _+^A(x)& \psi _{}^A(x)\\ \psi _{}^A(x)& \psi _+^A(x)\end{array}\right),\mathrm{\Psi }^B(x)=\left(\begin{array}{cc}\psi _+^B(x)& \psi _{}^B(x)\\ \psi _{}^B(x)& \psi _+^B(x)\end{array}\right).$$
(5.2)
It should be noted that $`\mathrm{\Psi }^{A,B}(x)`$ is not an independent field, but consists of the same Grassmann fields $`\psi _\pm ^{A,B}(x)`$ and $`\psi _\pm ^{A,B}(x)`$ as $`\mathrm{\Psi }^{A,B}(x)`$.
It should be pointed out that, since they emerge dynamically, the continuum fields of the low-energy effective theory can not be derived explicitly from the lattice operators of the microscopic Hubbard model. Still, the Grassmann fields $`\mathrm{\Psi }^{A,B}(x)`$ describing electrons and holes in the low-energy effective theory transform just like the lattice operators $`\mathrm{\Psi }_x^{A,B}`$ discussed before. In contrast to the lattice operators, the fields $`\mathrm{\Psi }^{A,B}(x)`$ are defined in the continuum. Hence, under the displacement symmetries $`D`$ and $`D^{}`$ one no longer distinguishes between the points $`x`$ and $`x+\widehat{i}`$. As a result, the transformation rules of the various symmetries take the form
$`SU(2)_s:`$ $`\mathrm{\Psi }^{A,B}(x)^{}=h(x)\mathrm{\Psi }^{A,B}(x),`$
$`SU(2)_Q:`$ $`{}_{}{}^{\stackrel{}{Q}}\mathrm{\Psi }_{}^{A,B}(x)=\mathrm{\Psi }^{A,B}(x)\mathrm{\Omega }^T,`$
$`D:`$ $`{}_{}{}^{D}\mathrm{\Psi }_{}^{A,B}(x)=\tau (x)\mathrm{\Psi }^{B,A}(x)\sigma _3,`$
$`D^{}:`$ $`{}_{}{}^{D^{}}\mathrm{\Psi }_{}^{A,B}(x)=(i\sigma _2)\mathrm{\Psi }^{B,A}(x)\sigma _3.`$ (5.3)
In components the symmetry transformations read
$`SU(2)_s:`$ $`\psi _\pm ^{A,B}(x)^{}=\mathrm{exp}(\pm i\alpha (x))\psi _\pm ^{A,B}(x),`$
$`U(1)_Q:`$ $`{}_{}{}^{Q}\psi _{\pm }^{A,B}(x)=\mathrm{exp}(i\omega )\psi _\pm ^{A,B}(x),`$
$`D:`$ $`{}_{}{}^{D}\psi _{\pm }^{A,B}(x)=\mathrm{exp}(i\phi (x))\psi _{}^{B,A}(x),`$
$`D^{}:`$ $`{}_{}{}^{D^{}}\psi _{\pm }^{A,B}(x)=\pm \psi _{}^{B,A}(x).`$ (5.4)
Under the space-time symmetries, i.e. under the 90 degrees rotation $`O`$, the reflection $`R`$, time-reversal $`T`$, and its unbroken variant $`T^{}`$ the fermion fields transform as
$`O:`$ $`{}_{}{}^{O}\mathrm{\Psi }_{}^{A,B}(x)=\mathrm{\Psi }^{A,B}(Ox),`$
$`R:`$ $`{}_{}{}^{R}\mathrm{\Psi }_{}^{A,B}(x)=\mathrm{\Psi }^{A,B}(Rx),`$
$`T:`$ $`{}_{}{}^{T}\mathrm{\Psi }_{}^{A,B}(x)=\tau (Tx)(i\sigma _2)\mathrm{\Psi }^{A,B}(Tx)^T\sigma _3,`$
$`{}_{}{}^{T}\mathrm{\Psi }_{}^{A,B}(x)=\sigma _3\mathrm{\Psi }^{A,B}(Tx)^T(i\sigma _2)^{}\tau (Tx)^{},`$
$`T^{}:`$ $`{}_{}{}^{T^{}}\mathrm{\Psi }_{}^{A,B}(x)=\mathrm{\Psi }^{A,B}(Tx)^T\sigma _3,`$ (5.5)
$`{}_{}{}^{T^{}}\mathrm{\Psi }_{}^{A,B}(x)=\sigma _3\mathrm{\Psi }^{A,B}(Tx)^T.`$
Again an upper index $`T`$ on the right denotes transpose, while on the left it denotes time-reversal. The form of the time-reversal symmetry $`T`$ in the effective theory with non-linearly realized $`SU(2)_s`$ symmetry follows from the usual form of time-reversal in the Euclidean path integral of a non-relativistic theory in which the spin symmetry is linearly realized. The fermion fields in the two formulations just differ by a factor $`u(x)`$. In components the previous relations take the form
$`O:`$ $`{}_{}{}^{O}\psi _{\pm }^{A,B}(x)=\psi _\pm ^{A,B}(Ox),`$
$`R:`$ $`{}_{}{}^{R}\psi _{\pm }^{A,B}(x)=\psi _\pm ^{A,B}(Rx),`$
$`T:`$ $`{}_{}{}^{T}\psi _{\pm }^{A,B}(x)=\mathrm{exp}(i\phi (Tx))\psi _\pm ^{A,B}(Tx),`$
$`{}_{}{}^{T}\psi _{\pm }^{A,B}(x)=\mathrm{exp}(\pm i\phi (Tx))\psi _\pm ^{A,B}(Tx),`$
$`T^{}:`$ $`{}_{}{}^{T^{}}\psi _{\pm }^{A,B}(x)=\psi _\pm ^{A,B}(Tx),`$ (5.6)
$`{}_{}{}^{T^{}}\psi _{\pm }^{A,B}(x)=\psi _\pm ^{A,B}(Tx).`$
It should be noted that the components $`+`$ and $``$ (denoting spin parallel and antiparallel to the direction of the staggered magnetization) are not interchanged under time-reversal. While both the spin of the fermion and the staggered magnetization change sign under time-reversal, the projection of one onto the other does not.
The action to be constructed in the next section must be invariant under the internal symmetries $`SU(2)_s`$, $`U(1)_Q`$ (or even $`SU(2)_Q`$), $`D`$ and $`D^{}`$, as well as under space-time translations and the other space-time symmetries $`O`$, $`R`$, and $`T`$ (or equivalently $`T^{}`$).
The fundamental forces underlying condensed matter physics are Poincarรฉ-invariant. However, some of the space-time symmetries may be spontaneously broken by the formation of a crystal lattice. The resulting Nambu-Goldstone bosons are the phonons, which play a central role in ordinary low-temperature superconductors by providing the force that binds Cooper pairs. In high-$`T_c`$ superconductors, on the other hand, it is expected that phonons alone cannot provide the mechanism for Cooper pair formation. In the Hubbard model (and also in our effective theory) phonons are explicitly excluded because one imposes a rigid lattice by hand. This does not only break continuous translations and rotations down to their discrete counterparts; it also breaks space-time rotations. In a relativistic context these would be the boosts of the Poincarรฉ group. In a non-relativistic theory the lattice explicitly breaks Galilean boost invariance, thus providing a preferred rest frame (a condensed matter โetherโ). As a consequence, the magnon-mediated forces between a pair of electrons or holes may depend on the center of mass momentum of the pair. In the actual high-$`T_c`$ materials Galilean (or more precisely Poincarรฉ symmetry) is spontaneously (and not explicitly) broken. If phonons play an important role in the understanding of high-temperature superconductivity, one should construct an effective theory of spontaneously broken (and thus non-linearly realized) $`SU(2)_s`$ and Galilean symmetry which would automatically include both magnons and phonons. This is indeed possible and presently under investigation using the techniques of low-energy effective field theory. In the present paper we assume that phonons play no major role in the cuprates. In that case, it is legitimate to break Galilean invariance explicitly instead of spontaneously.
### 5.2 Effective Action for Magnons and Charge Carriers
We now construct the leading terms in the effective action of magnons and electrons or holes. The effective theory provides a systematic low-energy expansion organized according to the number of derivatives in the terms of the effective action. We decompose the effective Lagrangian into an $`SU(2)_Q`$-invariant part $``$ and an $`SU(2)_Q`$-breaking (but still $`U(1)_Q`$-invariant) part $`\stackrel{~}{}`$. The contributions $`_{n_t,n_i,n_\psi }`$ and $`\stackrel{~}{}_{n_t,n_i,n_\psi }`$ to the effective Lagrangian are classified according to the number of time-derivatives $`n_t`$, the number of spatial derivatives $`n_i`$, and the number of fermion fields $`n_\psi `$ they contain. The total action is then given by
$$S[\psi _\pm ^{A,B},\psi _\pm ^{A,B},P]=d^2x๐t\underset{n_t,n_i,n_\psi }{}(_{n_t,n_i,n_\psi }+\stackrel{~}{}_{n_t,n_i,n_\psi })$$
(5.7)
and the partition function takes the form
$$Z=๐\psi _\pm ^{A,B}๐\psi _\pm ^{A,B}๐P\mathrm{exp}(S[\psi _\pm ^{A,B},\psi _\pm ^{A,B},P]).$$
(5.8)
Until now we have constructed the effective action in the $`Q=0`$ sector, i.e. for a half-filled system which is described entirely in terms of magnons. Since antiferromagnetic magnons have a โrelativisticโ dispersion relation (with the spin-wave velocity $`c`$ playing the role of the velocity of light), in pure magnon chiral perturbation theory one counts temporal and spatial derivatives as being of the same order. The leading contributions of eq.(3.5) take the form
$$_{2,0,0}=\frac{\rho _s}{c^2}\text{Tr}[_tP_tP],_{0,2,0}=\rho _s\text{Tr}[_iP_iP].$$
(5.9)
Next we consider terms quadratic in the fermion fields. These contribute to the scattering of magnons off electrons or holes in the $`Q=\pm 1`$ sectors and they generally describe the propagation of charge carriers in an antiferromagnet with $`|Q|1`$. In contrast to magnons, electrons or holes are massive and have a non-relativistic dispersion relation. Hence, it is natural to count one temporal and two spatial derivatives as being of the same order. In order to count derivatives consistently, in the $`Q0`$ sectors it may thus be necessary to also consider the pure magnon term $`_{2,0,0}`$ with two temporal derivatives as being of higher order. The leading order terms without any derivatives which are Hermitean and invariant under $`SU(2)_s`$, $`SU(2)_Q`$, $`D`$, and $`D^{}`$ as well as under the space-time symmetries $`O`$, $`R`$, $`T`$, and $`T^{}`$ take the form
$`_{0,0,2}`$ $`=`$ $`M_1\text{Tr}[\mathrm{\Psi }^A\mathrm{\Psi }^B]+{\displaystyle \frac{M_2}{2}}\text{Tr}[\mathrm{\Psi }^A\sigma _3\mathrm{\Psi }^A\mathrm{\Psi }^B\sigma _3\mathrm{\Psi }^B]`$ (5.10)
$`=`$ $`M_1(\psi _+^A\psi _+^B+\psi _{}^A\psi _{}^B+\psi _+^B\psi _+^A+\psi _{}^B\psi _{}^A)`$
$`+M_2(\psi _+^A\psi _+^A\psi _{}^A\psi _{}^A\psi _+^B\psi _+^B+\psi _{}^B\psi _{}^B).`$
The mass parameters $`M_1`$ and $`M_2`$ (as well as all other low-energy parameters to be introduced below) take real values in order to ensure Hermiticity of the corresponding Hamiltonian. It should be noted that
$`\text{Tr}[\mathrm{\Psi }^A\mathrm{\Psi }^A]=\text{Tr}[\mathrm{\Psi }^B\mathrm{\Psi }^B]=0,`$
$`\text{Tr}[\mathrm{\Psi }^A\mathrm{\Psi }^B]=\text{Tr}[\mathrm{\Psi }^B\mathrm{\Psi }^A],`$ (5.11)
due to the anticommutativity of the Grassmann fields. When we impose only the generic $`U(1)_Q`$ but not the full $`SU(2)_Q`$ symmetry, one more fermion mass term can be added
$`\stackrel{~}{}_{0,0,2}`$ $`=`$ $`{\displaystyle \frac{m}{2}}\text{Tr}[\mathrm{\Psi }^A\mathrm{\Psi }^A\sigma _3+\mathrm{\Psi }^B\mathrm{\Psi }^B\sigma _3]`$ (5.12)
$`=`$ $`m(\psi _+^A\psi _+^A+\psi _{}^A\psi _{}^A+\psi _+^B\psi _+^B+\psi _{}^B\psi _{}^B).`$
This term can be absorbed into a redefinition of the chemical potential. Remarkably, no other fermion mass terms (consistent with the $`SU(2)_s`$, $`U(1)_Q`$, $`D`$, $`D^{}`$, $`T`$, and $`T^{}`$ symmetries) exist. In particular, it is useful to note that
$$\text{Tr}[\mathrm{\Psi }^A\sigma _3\mathrm{\Psi }^A\sigma _3]=\text{Tr}[\mathrm{\Psi }^B\sigma _3\mathrm{\Psi }^B\sigma _3]=0.$$
(5.13)
The terms with one temporal derivative are given by
$`_{1,0,2}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\text{Tr}[\mathrm{\Psi }^AD_t\mathrm{\Psi }^A+\mathrm{\Psi }^BD_t\mathrm{\Psi }^B]`$ (5.14)
$`+{\displaystyle \frac{\mathrm{\Lambda }_1}{2}}\text{Tr}[\mathrm{\Psi }^AV_t\mathrm{\Psi }^A+\mathrm{\Psi }^BV_t\mathrm{\Psi }^B]+\mathrm{\Lambda }_2\text{Tr}[\mathrm{\Psi }^A\sigma _3V_t\mathrm{\Psi }^B]`$
$`=`$ $`\psi _+^AD_t\psi _+^A+\psi _{}^AD_t\psi _{}^A+\psi _+^BD_t\psi _+^B+\psi _{}^BD_t\psi _{}^B`$
$`+\mathrm{\Lambda }_1(\psi _+^Av_t^+\psi _{}^A+\psi _{}^Av_t^{}\psi _+^A+\psi _+^Bv_t^+\psi _{}^B+\psi _{}^Bv_t^{}\psi _+^B)`$
$`+\mathrm{\Lambda }_2(\psi _+^Av_t^+\psi _{}^B+\psi _{}^Bv_t^{}\psi _+^A\psi _+^Bv_t^+\psi _{}^A\psi _{}^Av_t^{}\psi _+^B).`$
Here $`V_t`$ is the field defined in eq.(3.45) and the covariant derivatives are those of eq.(4.2). In components they take the form
$`D_\mu \psi _\pm ^{A,B}(x)=(_\mu \pm iv_\mu ^3(x))\psi _\pm ^{A,B}(x),`$
$`D_\mu \psi _\pm ^{A,B}(x)=(_\mu iv_\mu ^3(x))\psi _\pm ^{A,B}(x).`$ (5.15)
Note that $`v_t^3`$ as well as $`v_t^\pm `$ (and hence $`V_t`$) count like one temporal derivative because these composite fields indeed contain one time-derivative of the magnon field.
When one derives the Euclidean path integral from the Hamiltonian formulation of the effective theory, the term $`\psi _+^A_t\psi _+^A+\psi _{}^A_t\psi _{}^A+\psi _+^B_t\psi _+^B+\psi _{}^B_t\psi _{}^B`$ arises from the pairs of anticommuting fermion operators. It should be noted that there are two more $`SU(2)_Q`$-breaking but $`U(1)_Q`$-invariant terms with a single time-derivative
$`{\displaystyle \frac{1}{2}}\text{Tr}[\mathrm{\Psi }^A\sigma _3D_t\mathrm{\Psi }^A\sigma _3\mathrm{\Psi }^B\sigma _3D_t\mathrm{\Psi }^B\sigma _3]`$
$`=\psi _+^AD_t\psi _+^A\psi _{}^AD_t\psi _{}^A\psi _+^BD_t\psi _+^B+\psi _{}^BD_t\psi _{}^B,`$
$`{\displaystyle \frac{1}{2}}\text{Tr}[\mathrm{\Psi }^AD_t\mathrm{\Psi }^B\sigma _3+\mathrm{\Psi }^BD_t\mathrm{\Psi }^A\sigma _3]`$
$`=\psi _+^AD_t\psi _+^B+\psi _{}^AD_t\psi _{}^B+\psi _+^BD_t\psi _+^A+\psi _{}^BD_t\psi _{}^A.`$ (5.16)
These terms need not be included in the effective Lagrangian, since they would not imply canonical anticommutation relations in the Hamiltonian formulation. In any case, as discussed in appendix B, if one does include these terms they can again be removed by an appropriate field redefinition.
Interestingly, there is only one more term that violates the $`SU(2)_Q`$ symmetry but still respects the $`U(1)_Q`$ symmetry
$`\stackrel{~}{}_{1,0,2}`$ $`=`$ $`\lambda \text{Tr}[\mathrm{\Psi }^AV_t\mathrm{\Psi }^B\sigma _3]`$ (5.17)
$`=`$ $`\lambda (\psi _+^Av_t^+\psi _{}^B+\psi _{}^Bv_t^{}\psi _+^A+\psi _+^Bv_t^+\psi _{}^A+\psi _{}^Av_t^{}\psi _+^B).`$
Further potential contributions are absent because, for example,
$`\text{Tr}[\mathrm{\Psi }^AD_t\mathrm{\Psi }^A\sigma _3]=\text{Tr}[\mathrm{\Psi }^BD_t\mathrm{\Psi }^B\sigma _3]=0,`$
$`\text{Tr}[\mathrm{\Psi }^AV_t\mathrm{\Psi }^A\sigma _3]=\text{Tr}[\mathrm{\Psi }^BV_t\mathrm{\Psi }^B\sigma _3]=0,`$
$`\text{Tr}[\mathrm{\Psi }^A\sigma _3V_t\mathrm{\Psi }^A\sigma _3]=\text{Tr}[\mathrm{\Psi }^B\sigma _3V_t\mathrm{\Psi }^B\sigma _3]=0.`$ (5.18)
Terms with a single spatial derivative are forbidden due to the reflection symmetry $`R`$ and the 90 degrees rotation symmetry $`O`$ of the quadratic spatial lattice of the underlying microscopic system. The terms with two spatial derivatives are given by
$`_{0,2,2}`$ $`=`$ $`{\displaystyle \frac{1}{2M_1^{}}}\text{Tr}[D_i\mathrm{\Psi }^AD_i\mathrm{\Psi }^B]+{\displaystyle \frac{1}{4M_2^{}}}\text{Tr}[D_i\mathrm{\Psi }^A\sigma _3D_i\mathrm{\Psi }^AD_i\mathrm{\Psi }^B\sigma _3D_i\mathrm{\Psi }^B]`$ (5.19)
$`+iK_1\text{Tr}[D_i\mathrm{\Psi }^AV_i\mathrm{\Psi }^B+D_i\mathrm{\Psi }^BV_i\mathrm{\Psi }^A]`$
$`+iK_2\text{Tr}[D_i\mathrm{\Psi }^A\sigma _3V_i\mathrm{\Psi }^AD_i\mathrm{\Psi }^B\sigma _3V_i\mathrm{\Psi }^B]`$
$`+N_1\text{Tr}[\mathrm{\Psi }^AV_iV_i\mathrm{\Psi }^B]+{\displaystyle \frac{N_2}{2}}\text{Tr}[\mathrm{\Psi }^AV_i\sigma _3V_i\mathrm{\Psi }^A\mathrm{\Psi }^BV_i\sigma _3V_i\mathrm{\Psi }^B]`$
$`=`$ $`{\displaystyle \frac{1}{2M_1^{}}}(D_i\psi _+^AD_i\psi _+^B+D_i\psi _+^BD_i\psi _+^A+D_i\psi _{}^AD_i\psi _{}^B+D_i\psi _{}^BD_i\psi _{}^A)`$
$`+{\displaystyle \frac{1}{2M_2^{}}}(D_i\psi _+^AD_i\psi _+^AD_i\psi _{}^AD_i\psi _{}^AD_i\psi _+^BD_i\psi _+^B+D_i\psi _{}^BD_i\psi _{}^B)`$
$`+iK_1(D_i\psi _+^Av_i^+\psi _{}^B\psi _{}^Bv_i^{}D_i\psi _+^A+D_i\psi _{}^Av_i^{}\psi _+^B\psi _+^Bv_i^+D_i\psi _{}^A`$
$`+D_i\psi _+^Bv_i^+\psi _{}^A\psi _{}^Av_i^{}D_i\psi _+^B+D_i\psi _{}^Bv_i^{}\psi _+^A\psi _+^Av_i^+D_i\psi _{}^B)`$
$`+iK_2(D_i\psi _+^Av_i^+\psi _{}^A\psi _{}^Av_i^{}D_i\psi _+^AD_i\psi _{}^Av_i^{}\psi _+^A+\psi _+^Av_i^+D_i\psi _{}^A`$
$`D_i\psi _+^Bv_i^+\psi _{}^B+\psi _{}^Bv_i^{}D_i\psi _+^B+D_i\psi _{}^Bv_i^{}\psi _+^B\psi _+^Bv_i^+D_i\psi _{}^B)`$
$`+N_1(\psi _+^Av_i^+v_i^{}\psi _+^B+\psi _{}^Av_i^{}v_i^+\psi _{}^B+\psi _+^Bv_i^+v_i^{}\psi _+^A+\psi _{}^Bv_i^{}v_i^+\psi _{}^A)`$
$`N_2(\psi _+^Av_i^+v_i^{}\psi _+^A\psi _{}^Av_i^{}v_i^+\psi _{}^A\psi _+^Bv_i^+v_i^{}\psi _+^B+\psi _{}^Bv_i^{}v_i^+\psi _{}^B).`$
Note that the imaginary unit $`i`$ in front of the terms proportional to $`K_1`$ and $`K_2`$ is necessary to ensure that the corresponding Hamiltonian is Hermitean. In principle, terms containing $`D_iD_i`$ and $`D_iV_i`$ could also be written down. However, upon partial integration, up to irrelevant surface terms they lead to the same Euclidean action as the terms constructed here. Since the doped electrons or holes are non-relativistic, there is no reason why the kinetic mass parameters $`M_1^{}`$ and $`M_2^{}`$ should agree with the rest mass parameters $`M_1`$ and $`M_2`$. In addition, there are again terms that break the $`SU(2)_Q`$ symmetry but leave the $`U(1)_Q`$ symmetry intact
$`\stackrel{~}{}_{0,2,2}`$ $`=`$ $`{\displaystyle \frac{1}{4m^{}}}\text{Tr}[D_i\mathrm{\Psi }^AD_i\mathrm{\Psi }^A\sigma _3+D_i\mathrm{\Psi }^BD_i\mathrm{\Psi }^B\sigma _3]`$ (5.20)
$`+i\kappa _1\text{Tr}[D_i\mathrm{\Psi }^A\sigma _3V_i\mathrm{\Psi }^B\sigma _3+D_i\mathrm{\Psi }^BV_i\sigma _3\mathrm{\Psi }^A\sigma _3]`$
$`+i\kappa _2\text{Tr}[D_i\mathrm{\Psi }^AV_i\mathrm{\Psi }^A\sigma _3+D_i\mathrm{\Psi }^BV_i\mathrm{\Psi }^B\sigma _3]`$
$`+{\displaystyle \frac{\nu }{2}}\text{Tr}[\mathrm{\Psi }^AV_iV_i\mathrm{\Psi }^A\sigma _3+\mathrm{\Psi }^BV_iV_i\mathrm{\Psi }^B\sigma _3]`$
$`=`$ $`{\displaystyle \frac{1}{2m^{}}}(D_i\psi _+^AD_i\psi _+^A+D_i\psi _{}^AD_i\psi _{}^A+D_i\psi _+^BD_i\psi _+^B+D_i\psi _{}^BD_i\psi _{}^B)`$
$`+i\kappa _1(D_i\psi _+^Av_i^+\psi _{}^B\psi _{}^Bv_i^{}D_i\psi _+^AD_i\psi _{}^Av_i^{}\psi _+^B+\psi _+^Bv_i^+D_i\psi _{}^A`$
$`D_i\psi _+^Bv_i^+\psi _{}^A+\psi _{}^Av_i^{}D_i\psi _+^B+D_i\psi _{}^Bv_i^{}\psi _+^A\psi _+^Av_i^+D_i\psi _{}^B)`$
$`+i\kappa _2(D_i\psi _+^Av_i^+\psi _{}^A\psi _{}^Av_i^{}D_i\psi _+^A+D_i\psi _{}^Av_i^{}\psi _+^A\psi _+^Av_i^+D_i\psi _{}^A`$
$`+D_i\psi _+^Bv_i^+\psi _{}^B\psi _{}^Bv_i^{}D_i\psi _+^B+D_i\psi _{}^Bv_i^{}\psi _+^B\psi _+^Bv_i^+D_i\psi _{}^B)`$
$`+\nu (\psi _+^Av_i^+v_i^{}\psi _+^A+\psi _{}^Av_i^{}v_i^+\psi _{}^A+\psi _+^Bv_i^+v_i^{}\psi _+^B+\psi _{}^Bv_i^{}v_i^+\psi _{}^B).`$
Next we consider terms quartic in the fermion fields which describe short-range interactions between the charge carriers. To lowest order there are five linearly independent 4-fermion contact interaction terms
$`_{0,0,4}`$ $`=`$ $`{\displaystyle \frac{G_1}{12}}\text{Tr}[\mathrm{\Psi }^A\mathrm{\Psi }^A\mathrm{\Psi }^A\mathrm{\Psi }^A+\mathrm{\Psi }^B\mathrm{\Psi }^B\mathrm{\Psi }^B\mathrm{\Psi }^B]+{\displaystyle \frac{G_2}{2}}\text{Tr}[\mathrm{\Psi }^A\mathrm{\Psi }^A\mathrm{\Psi }^B\mathrm{\Psi }^B]`$
$`+{\displaystyle \frac{G_3}{2}}\text{Tr}[\mathrm{\Psi }^A\mathrm{\Psi }^B\mathrm{\Psi }^B\mathrm{\Psi }^A]+{\displaystyle \frac{G_4}{2}}\text{Tr}[\mathrm{\Psi }^A\sigma _3\mathrm{\Psi }^A\mathrm{\Psi }^B\sigma _3\mathrm{\Psi }^B]`$
$`+G_5\text{Tr}[\mathrm{\Psi }^A\sigma _3\mathrm{\Psi }^A\mathrm{\Psi }^A\mathrm{\Psi }^B\mathrm{\Psi }^B\sigma _3\mathrm{\Psi }^B\mathrm{\Psi }^B\mathrm{\Psi }^A]`$
$`=`$ $`G_1(\psi _+^A\psi _+^A\psi _{}^A\psi _{}^A+\psi _+^B\psi _+^B\psi _{}^B\psi _{}^B)`$
$`+G_2(\psi _+^A\psi _+^A\psi _+^B\psi _+^B+\psi _+^A\psi _+^A\psi _{}^B\psi _{}^B+\psi _{}^A\psi _{}^A\psi _+^B\psi _+^B+\psi _{}^A\psi _{}^A\psi _{}^B\psi _{}^B`$
$`2\psi _+^A\psi _+^B\psi _{}^A\psi _{}^B2\psi _+^B\psi _+^A\psi _{}^B\psi _{}^A)`$
$`+G_3(\psi _+^A\psi _+^A\psi _{}^B\psi _{}^B+\psi _{}^A\psi _{}^A\psi _+^B\psi _+^B\psi _+^A\psi _+^A\psi _+^B\psi _+^B\psi _{}^A\psi _{}^A\psi _{}^B\psi _{}^B`$
$`2\psi _+^A\psi _{}^A\psi _{}^B\psi _+^B2\psi _{}^A\psi _+^A\psi _+^B\psi _{}^B)`$
$`+G_4(\psi _+^A\psi _+^A\psi _+^B\psi _+^B\psi _+^A\psi _+^A\psi _{}^B\psi _{}^B\psi _{}^A\psi _{}^A\psi _+^B\psi _+^B+\psi _{}^A\psi _{}^A\psi _{}^B\psi _{}^B)`$
$`+G_5(\psi _+^A\psi _+^A\psi _{}^A\psi _{}^B+\psi _+^A\psi _+^A\psi _{}^B\psi _{}^A\psi _+^B\psi _+^B\psi _{}^A\psi _{}^B\psi _+^B\psi _+^B\psi _{}^B\psi _{}^A`$
$`\psi _+^A\psi _+^B\psi _{}^A\psi _{}^A\psi _+^B\psi _+^A\psi _{}^A\psi _{}^A+\psi _+^A\psi _+^B\psi _{}^B\psi _{}^B+\psi _+^B\psi _+^A\psi _{}^B\psi _{}^B).`$
It is interesting to note that
$`\text{Tr}[\mathrm{\Psi }^A\mathrm{\Psi }^B\mathrm{\Psi }^B\mathrm{\Psi }^A+\mathrm{\Psi }^A\mathrm{\Psi }^A\mathrm{\Psi }^B\mathrm{\Psi }^B+2\mathrm{\Psi }^A\mathrm{\Psi }^B\mathrm{\Psi }^A\mathrm{\Psi }^B]=0,`$
$`\text{Tr}[\mathrm{\Psi }^A\mathrm{\Psi }^B\mathrm{\Psi }^B\mathrm{\Psi }^A+\mathrm{\Psi }^A\sigma _3\mathrm{\Psi }^B\mathrm{\Psi }^B\sigma _3\mathrm{\Psi }^A+2\mathrm{\Psi }^A\sigma _3\mathrm{\Psi }^A\mathrm{\Psi }^B\sigma _3\mathrm{\Psi }^B]=0,`$
$`\text{Tr}[\mathrm{\Psi }^A\mathrm{\Psi }^B\mathrm{\Psi }^B\mathrm{\Psi }^A\mathrm{\Psi }^A\mathrm{\Psi }^A\mathrm{\Psi }^B\mathrm{\Psi }^B]=2(\text{Tr}[\mathrm{\Psi }^A\mathrm{\Psi }^B])^2.`$ (5.22)
Together with other relations similar to these ones, this implies that the terms listed above form a maximal linearly independent set.
Again, there are additional terms that are invariant under $`U(1)_Q`$ but not under $`SU(2)_Q`$
$`\stackrel{~}{}_{0,0,4}`$ $`=`$ $`{\displaystyle \frac{g_1}{4}}\text{Tr}[\mathrm{\Psi }^A\sigma _3\mathrm{\Psi }^A\mathrm{\Psi }^B\mathrm{\Psi }^B\sigma _3\mathrm{\Psi }^B\sigma _3\mathrm{\Psi }^B\mathrm{\Psi }^A\mathrm{\Psi }^A\sigma _3]+{\displaystyle \frac{g_2}{2}}\text{Tr}[\mathrm{\Psi }^A\mathrm{\Psi }^A\sigma _3\mathrm{\Psi }^B\mathrm{\Psi }^B\sigma _3]`$
$`+g_3\text{Tr}[\mathrm{\Psi }^A\mathrm{\Psi }^A\sigma _3\mathrm{\Psi }^A\mathrm{\Psi }^B+\mathrm{\Psi }^B\mathrm{\Psi }^B\sigma _3\mathrm{\Psi }^B\mathrm{\Psi }^A]`$
$`=`$ $`g_1(\psi _+^A\psi _+^A\psi _{}^B\psi _{}^B\psi _{}^A\psi _{}^A\psi _+^B\psi _+^B)`$
$`+g_2(\psi _+^A\psi _+^A\psi _+^B\psi _+^B+\psi _+^A\psi _+^A\psi _{}^B\psi _{}^B+\psi _{}^A\psi _{}^A\psi _+^B\psi _+^B+\psi _{}^A\psi _{}^A\psi _{}^B\psi _{}^B`$
$`+2\psi _+^A\psi _+^B\psi _{}^A\psi _{}^B+2\psi _+^B\psi _+^A\psi _{}^B\psi _{}^A)`$
$`g_3(\psi _+^A\psi _+^B\psi _{}^A\psi _{}^A+\psi _+^A\psi _+^B\psi _{}^B\psi _{}^B+\psi _{}^A\psi _{}^B\psi _+^A\psi _+^A+\psi _{}^A\psi _{}^B\psi _+^B\psi _+^B`$
$`+\psi _+^B\psi _+^A\psi _{}^A\psi _{}^A+\psi _+^B\psi _+^A\psi _{}^B\psi _{}^B+\psi _{}^B\psi _{}^A\psi _+^A\psi _+^A+\psi _{}^B\psi _{}^A\psi _+^B\psi _+^B).`$
One may note that, for example,
$$\text{Tr}[\mathrm{\Psi }^A\mathrm{\Psi }^A\mathrm{\Psi }^B\mathrm{\Psi }^B+\mathrm{\Psi }^A\mathrm{\Psi }^A\sigma _3\mathrm{\Psi }^B\mathrm{\Psi }^B\sigma _3+2\mathrm{\Psi }^A\mathrm{\Psi }^B\sigma _3\mathrm{\Psi }^B\mathrm{\Psi }^A\sigma _3]=0.$$
(5.24)
Together with further relations of a similar kind, this implies that there are no other linearly independent 4-fermion terms that obey the relevant symmetries.
For completeness, let us also construct the terms containing six fermion fields and no derivatives. The $`SU(2)_Q`$-invariant 6-fermion terms can be written as
$`_{0,0,6}`$ $`=`$ $`{\displaystyle \frac{H_1}{4}}\text{Tr}[\mathrm{\Psi }^A\sigma _3\mathrm{\Psi }^A\mathrm{\Psi }^A\sigma _3\mathrm{\Psi }^A\mathrm{\Psi }^B\sigma _3\mathrm{\Psi }^B\mathrm{\Psi }^B\sigma _3\mathrm{\Psi }^B\mathrm{\Psi }^B\sigma _3\mathrm{\Psi }^B\mathrm{\Psi }^A\sigma _3\mathrm{\Psi }^A]`$ (5.25)
$`+{\displaystyle \frac{H_2}{3}}\text{Tr}[\mathrm{\Psi }^A\mathrm{\Psi }^B\mathrm{\Psi }^A\mathrm{\Psi }^B\mathrm{\Psi }^A\mathrm{\Psi }^B]`$
$`=`$ $`H_1(\psi _+^A\psi _+^A\psi _+^B\psi _+^B\psi _{}^B\psi _{}^B+\psi _+^A\psi _+^A\psi _{}^A\psi _{}^A\psi _{}^B\psi _{}^B`$
$`\psi _{}^A\psi _{}^A\psi _+^B\psi _+^B\psi _{}^B\psi _{}^B\psi _+^A\psi _+^A\psi _{}^A\psi _{}^A\psi _+^B\psi _+^B)`$
$`+H_2(\psi _+^A\psi _+^A\psi _+^B\psi _+^B\psi _{}^B\psi _{}^A+\psi _{}^A\psi _{}^A\psi _{}^B\psi _{}^B\psi _+^B\psi _+^A`$
$`+\psi _{}^A\psi _{}^A\psi _{}^B\psi _{}^B\psi _+^A\psi _+^B+\psi _+^A\psi _+^A\psi _+^B\psi _+^B\psi _{}^A\psi _{}^B).`$
It is interesting to note that
$`\text{Tr}[\mathrm{\Psi }^A\sigma _3\mathrm{\Psi }^A\mathrm{\Psi }^A\sigma _3\mathrm{\Psi }^A\mathrm{\Psi }^B\sigma _3\mathrm{\Psi }^B\mathrm{\Psi }^B\sigma _3\mathrm{\Psi }^B\mathrm{\Psi }^B\sigma _3\mathrm{\Psi }^B\mathrm{\Psi }^A\sigma _3\mathrm{\Psi }^A]`$
$`={\displaystyle \frac{1}{2}}\text{Tr}[\mathrm{\Psi }^A\sigma _3\mathrm{\Psi }^A\mathrm{\Psi }^B\sigma _3\mathrm{\Psi }^B]\text{Tr}[\mathrm{\Psi }^A\sigma _3\mathrm{\Psi }^A\mathrm{\Psi }^B\sigma _3\mathrm{\Psi }^B].`$ (5.26)
In addition, there is one $`SU(2)_Q`$-breaking (but $`U(1)_Q`$-invariant) 6-fermion term
$`\stackrel{~}{}_{0,0,6}`$ $`=`$ $`{\displaystyle \frac{h}{4}}\text{Tr}[\mathrm{\Psi }^A\mathrm{\Psi }^A\sigma _3\mathrm{\Psi }^A\mathrm{\Psi }^A\sigma _3\mathrm{\Psi }^B\mathrm{\Psi }^B\sigma _3+\mathrm{\Psi }^B\mathrm{\Psi }^B\sigma _3\mathrm{\Psi }^B\mathrm{\Psi }^B\sigma _3\mathrm{\Psi }^A\mathrm{\Psi }^A\sigma _3]`$ (5.27)
$`=`$ $`h(\psi _+^A\psi _+^A\psi _{}^A\psi _{}^A\psi _+^B\psi _+^B+\psi _+^A\psi _+^A\psi _{}^A\psi _{}^A\psi _{}^B\psi _{}^B`$
$`+\psi _+^A\psi _+^A\psi _+^B\psi _+^B\psi _{}^B\psi _{}^B+\psi _{}^A\psi _{}^A\psi _+^B\psi _+^B\psi _{}^B\psi _{}^B).`$
Finally, the only 8-fermion term with no derivatives takes the form
$`_{0,0,8}`$ $`=`$ $`{\displaystyle \frac{I}{24}}\text{Tr}[\mathrm{\Psi }^A\mathrm{\Psi }^A\mathrm{\Psi }^B\mathrm{\Psi }^B\mathrm{\Psi }^A\mathrm{\Psi }^A\mathrm{\Psi }^B\mathrm{\Psi }^B]`$ (5.28)
$`=`$ $`I\psi _+^A\psi _+^A\psi _+^B\psi _+^B\psi _{}^A\psi _{}^A\psi _{}^B\psi _{}^B,`$
which is $`SU(2)_Q`$-invariant. It may be noted that
$$\text{Tr}[\mathrm{\Psi }^A\mathrm{\Psi }^A\mathrm{\Psi }^B\mathrm{\Psi }^B\mathrm{\Psi }^A\mathrm{\Psi }^A\mathrm{\Psi }^B\mathrm{\Psi }^B]+\frac{1}{2}(\text{Tr}[\mathrm{\Psi }^A\mathrm{\Psi }^A\mathrm{\Psi }^B\mathrm{\Psi }^B])^2=0.$$
(5.29)
No $`SU(2)_Q`$-breaking 8-fermion term without derivatives exists, such that $`\stackrel{~}{}_{0,0,8}=0`$. Terms with more than eight fermion fields vanish due to the Pauli principle, unless one includes derivatives. Since such terms are of higher order than those without derivatives, they will not be constructed here. When one wants to address questions for which the short-distance forces between charge carriers are essential, it will be necessary to consider such terms. While constructing them is a straightforward exercise, it is not very illuminating and will hence be omitted at this stage.
The 4-, 6-, and 8-fermion contact terms parameterize short distance interactions with a large number of undetermined low-energy constants. Since this limits the predictive power of the effective theory at short distances, it is natural to concentrate on long-distance forces between the charge carriers. For example, one-magnon exchange mediates a long-range force that is unambiguously predicted by the effective theory in terms of just a few low-energy parameters.
It should be mentioned that there are many equivalent ways of rewriting the various contributions to the action in terms of traces. Hence, the above choices of terms are to some extent arbitrary. It is important that the selected terms form a maximal linearly independent set. For example, all determinants or products of traces of fermion fields can be written as linear combinations of the traces listed above. To verify the completeness and linear independence of the selected terms is a non-trivial task which was addressed by extensive use of the algebraic manipulation program FORM .
It is straightforward to include the fermion chemical potential $`\mu `$ in the effective theory. It appears as the temporal component of a purely imaginary $`U(1)_Q`$ gauge field and thus manifests itself in an additional contribution to the covariant derivative
$$D_t\mathrm{\Psi }^{A,B}(x)=_t\mathrm{\Psi }^{A,B}(x)+iv_t^3(x)\sigma _3\mathrm{\Psi }^{A,B}(x)\mu \mathrm{\Psi }^{A,B}(x)\sigma _3.$$
(5.30)
Before one can do consistent loop-calculations in the low-energy effective theory at non-zero $`Q`$ or at non-zero $`\mu `$ one must develop a power-counting scheme, e.g. along the lines of . This will be the subject of a future publication.
### 5.3 Dispersion Relations of Electrons and Holes
In this subsection, as an application of the effective theory, we consider the dispersion relations of the charge carriers. For this purpose we switch off the magnon field (i.e. $`P(x)=\frac{1}{2}(\mathrm{๐ฃ}\mathrm{๐ฃ}+\sigma _3)u(x)=\mathrm{๐ฃ}\mathrm{๐ฃ},v_\mu (x)=0`$) and consider the propagation of free charge carriers in the antiferromagnetic medium. In the absence of $`SU(2)_Q`$-breaking terms, the Lagrangian (quadratic in the fermion fields) then reduces to
$``$ $`=`$ $`\psi _+^A_t\psi _+^A+\psi _{}^A_t\psi _{}^A+\psi _+^B_t\psi _+^B+\psi _{}^B_t\psi _{}^B`$ (5.49)
$`+(\psi _+^A,\psi _+^B)\left(\begin{array}{cc}M_2& M_1\\ M_1& M_2\end{array}\right)\left(\begin{array}{c}\psi _+^A\\ \psi _+^B\end{array}\right)+(\psi _{}^A,\psi _{}^B)\left(\begin{array}{cc}M_2& M_1\\ M_1& M_2\end{array}\right)\left(\begin{array}{c}\psi _{}^A\\ \psi _{}^B\end{array}\right)`$
$`+(_i\psi _+^A,_i\psi _+^B)\left(\begin{array}{cc}\frac{1}{2M_2^{}}& \frac{1}{2M_1^{}}\\ \frac{1}{2M_1^{}}& \frac{1}{2M_2^{}}\end{array}\right)\left(\begin{array}{c}_i\psi _+^A\\ _i\psi _+^B\end{array}\right)`$
$`+(_i\psi _{}^A,_i\psi _{}^B)\left(\begin{array}{cc}\frac{1}{2M_2^{}}& \frac{1}{2M_1^{}}\\ \frac{1}{2M_1^{}}& \frac{1}{2M_2^{}}\end{array}\right)\left(\begin{array}{c}_i\psi _{}^A\\ _i\psi _{}^B\end{array}\right).`$
The eigenstates of free particles propagating with a 2-d momentum vector $`\stackrel{}{p}`$ arise as the eigenvectors of the matrices
$`H_+(p^2)`$ $`=`$ $`\left(\begin{array}{cc}M_2+\frac{p^2}{2M_2^{}}& M_1+\frac{p^2}{2M_1^{}}\\ M_1+\frac{p^2}{2M_1^{}}& M_2\frac{p^2}{2M_2^{}}\end{array}\right),`$ (5.52)
$`H_{}(p^2)`$ $`=`$ $`\left(\begin{array}{cc}M_2\frac{p^2}{2M_2^{}}& M_1+\frac{p^2}{2M_1^{}}\\ M_1+\frac{p^2}{2M_1^{}}& M_2+\frac{p^2}{2M_2^{}}\end{array}\right).`$ (5.55)
Due to the lack of Galilean invariance the eigenvectors depend on $`p^2`$, i.e. the probability for an electron or hole to be found on the $`A`$ or $`B`$ sublattice depends on the momentum. As a consequence of the displacement symmetries $`D`$ and $`D^{}`$ the eigenvalues of $`H_+(p^2)`$ and $`H_{}(p^2)`$ are the same. Both matrices have two eigenvalues
$`E_{1,2}(p^2)`$ $`=`$ $`\pm \sqrt{\left(M_1+{\displaystyle \frac{p^2}{2M_1^{}}}\right)^2+\left(M_2+{\displaystyle \frac{p^2}{2M_2^{}}}\right)^2}`$ (5.56)
$`=`$ $`\pm \left(M+{\displaystyle \frac{p^2}{2M^{}}}+๐ช(p^4)\right).`$
The positive energy states correspond to electrons, while the negative energy states correspond to holes. Not surprisingly, due to the $`SU(2)_Q`$ symmetry electrons and holes have the same dispersion relation. The rest mass $`M`$ and the kinetic mass $`M^{}`$ are given by
$$M=\sqrt{M_1^2+M_2^2},\frac{M}{M^{}}=\frac{M_1}{M_1^{}}+\frac{M_2}{M_2^{}}.$$
(5.57)
Next we take into account the additional terms that reduce the $`SU(2)_Q`$ symmetry to the $`U(1)_Q`$ symmetry. Then there are additional contributions to the energy
$$\stackrel{~}{H}_+(p^2)=\stackrel{~}{H}_{}(p^2)=\left(\begin{array}{cc}m+\frac{p^2}{2m^{}}& 0\\ 0& m+\frac{p^2}{2m^{}}\end{array}\right).$$
(5.58)
and the corresponding eigenvalues now take the form
$$E_{1,2}(p^2)=m+\frac{p^2}{2m^{}}\pm \left(M+\frac{p^2}{2M^{}}\right)+๐ช(p^4).$$
(5.59)
Still, the energies in the $`+`$ and $``$ sectors are the same. However, the electron and hole dispersion relations now differ.
At this point, we have constructed eigenstates of the free Hamiltonian with definite continuum momentum and with definite spin projection on the direction of the staggered magnetization. However, unlike the eigenstates of the underlying microscopic Hamiltonian, the states of the effective theory do not have a definite lattice momentum. Still, the low-energy effective theory defined in the continuum knows about the underlying lattice structure through the realization of the displacement symmetries $`D`$ and $`D^{}`$. Since the symmetry $`D`$ is spontaneously broken, neither the vacuum nor the single particle states are eigenstates of $`D`$. Operating twice with $`D`$ acts trivially on the fields, i.e. $`{}_{}{}^{DD}P(x)=P(x)`$, $`{}_{}{}^{DD}\mathrm{\Psi }_{\pm }^{A,B}=\mathrm{\Psi }_\pm ^{A,B}`$, and hence does not reveal any useful information. It is more useful to operate with the unbroken displacement symmetry $`D^{}`$. In particular, the vacuum state $`P(x)=\frac{1}{2}(\mathrm{๐ฃ}\mathrm{๐ฃ}+\sigma _3)`$ is invariant under $`D^{}`$. Still, in the way we constructed them, the electron or hole states of the effective theory are not eigenstates of $`D^{}`$. However, since states with spin parallel and antiparallel to the staggered magnetization are degenerate with each other, one can form appropriate linear combinations that are eigenstates of the displacement symmetry $`D^{}`$. Applying $`D^{}`$ twice one obtains
$${}_{}{}^{D^{}D^{}}\psi _{\pm }^{A,B}(x)=\pm ^D^{}\psi _{}^{B,A}(x)=\psi _\pm ^{A,B}(x),$$
(5.60)
which implies that the corresponding eigenvalue $`\lambda =\mathrm{exp}(ika)`$ of $`D^{}`$ obeys
$$\lambda ^2=\mathrm{exp}(2ika)=1ka=\pm \frac{\pi }{2}.$$
(5.61)
This is reminiscent of the result, mentioned in the introduction, that low-energy hole states are located at lattice momenta $`(\pm \frac{\pi }{2},\pm \frac{\pi }{2})`$ . However, the comparison with these findings is subtle. In particular, the results of the exact diagonalization study on small and of the Monte Carlo study on larger volumes must be interpreted carefully. In a finite volume (with periodic boundary conditions), in analogy to QCD , both the $`SU(2)_s`$ spin symmetry and the displacement symmetry $`D`$ are restored and the staggered magnetization acts as a quantum rotor . As a result, in contrast to the infinite volume limit, the single particle states in a finite volume can be constructed as eigenstates of $`D`$. It is interesting to note that the finite volume effects that lead to the restoration of the spontaneously broken symmetries $`SU(2)_s`$ and $`D`$ can be understood in the framework of the effective theory. This requires a nonperturbative quantum mechanical treatment along the lines of .
## 6 Systems with Holes only
In this section we consider the $`t`$-$`J`$ model as well as its low-energy effective theory. In the $`t`$-$`J`$ model holes are the only charge carriers which leads to substantial simplifications in the effective theory.
### 6.1 The $`t`$-$`J`$ Model
The $`t`$-$`J`$ model is defined by the Hamilton operator
$$H=P\left\{t\underset{x,i}{}(c_x^{}c_{x+\widehat{i}}+c_{x+\widehat{i}}^{}c_x)+J\underset{x,i}{}\stackrel{}{S}_x\stackrel{}{S}_{x+\widehat{i}}\mu \underset{x}{}(n_x1)\right\}P,$$
(6.1)
with
$$c_x=\left(\begin{array}{c}c_x\\ c_x\end{array}\right),S_x=c_x^{}\frac{\stackrel{}{\sigma }}{2}c_x,n_x=c_x^{}c_x.$$
(6.2)
In contrast to the Hubbard model, in the $`t`$-$`J`$ model the operators act in a restricted Hilbert space of empty or at most singly occupied sites. In particular, states with doubly occupied sites are exiled from the physical Hilbert space by the projection operator $`P`$. Hence, by definition, the $`t`$-$`J`$ model does not allow the addition of electrons to a half-filled state. Consequently, the only charge carriers are holes.
It is straightforward to show that the $`t`$-$`J`$ model has the same symmetries as the Hubbard model. The only exception is the $`SU(2)_Q`$ symmetry which relates electrons to holes in the Hubbard model, and which is absent in the $`t`$-$`J`$ model. Still, the Abelian fermion number symmetry $`U(1)_Q`$ remains exact in the $`t`$-$`J`$ model.
### 6.2 Effective Theory for Magnons and Holes
Since, up to the $`SU(2)_Q`$ symmetry, the $`t`$-$`J`$ model has the same symmetries as the Hubbard model, the effective theory of the previous section also applies in this case. Of course, the values of the low-energy parameters will be different than for the Hubbard model. Still, the absence of electrons beyond half-filling leads to drastic simplifications. In particular, in the effective theory the absence of electrons manifests itself by an infinite electron rest mass. Consequently, with a finite amount of energy these excitations cannot be generated. As discussed in the previous section, the diagonalization of the mass matrices of electrons and holes yields
$`U_\pm \left(\begin{array}{cc}m\pm M_2& M_1\\ M_1& mM_2\end{array}\right)U_\pm ^{}=\left(\begin{array}{cc}m\pm \sqrt{M_1^2+M_2^2}& 0\\ 0& m\sqrt{M_1^2+M_2^2}\end{array}\right),`$ (6.7)
$`U_\pm =\left(\begin{array}{cc}X& \pm Y\\ Y& X\end{array}\right),X,Y๐จ๐ฑ.`$ (6.10)
The eigenvectors corresponding to the eigenvalue $`m+\sqrt{M_1^2+M_2^2}`$ describe electrons, while the ones corresponding to $`m\sqrt{M_1^2+M_2^2}`$ describe holes. When the electron rest mass $`m+\sqrt{M_1^2+M_2^2}`$ goes to infinity, the corresponding eigenvector fields
$$X\psi _+^A(x)+Y\psi _+^B(x)=0,Y\psi _{}^A(x)+X\psi _{}^B(x)=0,$$
(6.11)
which describe electrons, must be put to zero. The orthogonal combinations
$$\psi _+(x)=Y\psi _+^A(x)+X\psi _+^B(x),\psi _{}(x)=X\psi _{}^A(x)Y\psi _{}^B(x),$$
(6.12)
describe holes and must be kept. As a result, the number of degrees of freedom is reduced by a factor of two. In complete analogy to the discussion in appendix B one can show that the hole field $`\psi _\pm (x)`$ transforms as follows under the various symmetry operations
$`SU(2)_s:`$ $`\psi _\pm (x)^{}=\mathrm{exp}(\pm i\alpha (x))\psi _\pm (x),`$
$`U(1)_Q:`$ $`{}_{}{}^{Q}\psi _{\pm }^{}(x)=\mathrm{exp}(i\omega )\psi _\pm (x),`$
$`D:`$ $`{}_{}{}^{D}\psi _{\pm }^{}(x)=\mathrm{exp}(i\phi (x))\psi _{}(x),`$
$`D^{}:`$ $`{}_{}{}^{D^{}}\psi _{\pm }^{}(x)=\pm \psi _{}(x),`$
$`O:`$ $`{}_{}{}^{O}\psi _{\pm }^{}(x)=\psi _\pm (Ox),`$
$`R:`$ $`{}_{}{}^{R}\psi _{\pm }^{}(x)=\psi _\pm (Rx),`$
$`T:`$ $`{}_{}{}^{T}\psi _{\pm }^{}(x)=\mathrm{exp}(i\phi (Tx))\psi _\pm ^{}(Tx),`$
$`{}_{}{}^{T}\psi _{\pm }^{}(x)=\mathrm{exp}(\pm i\phi (Tx))\psi _\pm (Tx),`$
$`T^{}:`$ $`{}_{}{}^{T^{}}\psi _{\pm }^{}(x)=\psi _\pm ^{}(Tx),`$ (6.13)
$`{}_{}{}^{T^{}}\psi _{\pm }^{}(x)=\psi _\pm (Tx).`$
Hence, except for the $`SU(2)_Q`$ symmetry, all symmetries can also be implemented on the hole fields alone. It should be noted that the transformation laws for $`\psi _\pm (x)`$ result from those for $`\psi _\pm ^{A,B}(x)`$ simply by dropping the sublattice indices $`A`$ and $`B`$.
The absence of electron fields also drastically reduces the number of terms one can write down in the low-energy effective theory. In particular, the leading terms in the effective action now take the form
$`S[\psi _\pm ^{},\psi _\pm ,P]`$ $`=`$ $`{\displaystyle }d^2xdt\{\rho _s\text{Tr}[_iP_iP+{\displaystyle \frac{1}{c^2}}_tP_tP]+M(\psi _+^{}\psi _++\psi _{}^{}\psi _{})`$ (6.14)
$`+\psi _+^{}D_t\psi _++\psi _{}^{}D_t\psi _{}+{\displaystyle \frac{1}{2M^{}}}(D_i\psi _+^{}D_i\psi _++D_i\psi _{}^{}D_i\psi _{})`$
$`+\mathrm{\Lambda }(\psi _+^{}v_t^+\psi _{}+\psi _{}^{}v_t^{}\psi _+)`$
$`+iK(D_i\psi _+^{}v_i^+\psi _{}\psi _{}^{}v_i^{}D_i\psi _++D_i\psi _{}^{}v_i^{}\psi _+\psi _+^{}v_i^+D_i\psi _{})`$
$`+N(\psi _+^{}v_i^+v_i^{}\psi _++\psi _{}^{}v_i^{}v_i^+\psi _{})+G\psi _+^{}\psi _+\psi _{}^{}\psi _{}\}.`$
This form of the effective action is similar to (but not identical with) the ones of . In particular, in some of those works spin-charge separation was invoked and spinless fermions were considered. Also the role of the sublattice indices (which have at this stage disappeared from our description) is different in those approaches. Furthermore, the dynamical role attributed to the composite gauge field in some of those works is different than in our effective theory. It should be pointed out that the above effective Lagrangian correctly describes the low-energy dynamics of holes only if electrons are completely absent beyond half-filling (as it is indeed the case in the $`t`$-$`J`$ model). Otherwise the general effective theory of the previous section with a larger number of low-energy constants (and thus with somewhat reduced predictive power) must be employed.
## 7 Coupling to External Electromagnetic Fields
In the following sections we will couple both microscopic and effective theories for antiferromagnets to external electromagnetic fields. For this purpose, we will make use of an observation by Frรถhlich and Studer concerning the Pauli equation .
### 7.1 Local $`SU(2)_s`$ Symmetry of the Pauli Equation
Up to corrections of order $`1/M_e^3`$ (where $`M_e`$ is the electron mass) the Pauli equation (i.e. the non-relativistic reduction of the Dirac equation to its upper components) takes the form
$$i(_tie\mathrm{\Phi }+i\frac{e}{8M_e^2}\stackrel{}{}\stackrel{}{E}+i\frac{e}{2M_e}\stackrel{}{B}\stackrel{}{\sigma })\mathrm{\Psi }=\frac{1}{2M_e}(\stackrel{}{}+ie\stackrel{}{A}i\frac{e}{4M_e}\stackrel{}{E}\times \stackrel{}{\sigma })^2\mathrm{\Psi }.$$
(7.1)
Here $`\mathrm{\Psi }(x)`$ is a 2-component Pauli spinor at the space-time point $`x=(\stackrel{}{x},t)`$, $`\stackrel{}{\sigma }`$ are the Pauli matrices, $`\mathrm{\Phi }(x)`$ and $`\stackrel{}{A}(x)`$ are the electromagnetic scalar and vector potentials, and
$$\stackrel{}{E}(x)=\stackrel{}{}\mathrm{\Phi }(x)_t\stackrel{}{A}(x),\stackrel{}{B}(x)=\stackrel{}{}\times \stackrel{}{A}(x),$$
(7.2)
are the usual electromagnetic field strengths. The first two terms on the left-hand side of eq.(7.1) form the $`U(1)_Q`$ covariant derivative familiar from QED. The third (Darwin) and fourth (Zeeman) term on the left-hand side represent relativistic corrections. The first two terms on the right-hand side again form an ordinary $`U(1)_Q`$ covariant derivative, while the third term represents the relativistic spin-orbit coupling. The Pauli equation transforms covariantly under $`U(1)_Q`$ gauge transformations
$${}_{}{}^{Q}\mathrm{\Psi }(x)=\mathrm{exp}(i\omega (x))\mathrm{\Psi }(x),^Q\mathrm{\Phi }(x)=\mathrm{\Phi }(x)+\frac{1}{e}_t\omega (x),^Q\stackrel{}{A}(x)=\stackrel{}{A}(x)\frac{1}{e}\stackrel{}{}\omega (x).$$
(7.3)
Obviously, it is also covariant under global spatial rotations
$${}_{}{}^{๐ช}\mathrm{\Psi }(\stackrel{}{x},t)=g\mathrm{\Psi }(๐ช\stackrel{}{x},t),^๐ช\mathrm{\Phi }(\stackrel{}{x},t)=\mathrm{\Phi }(๐ช\stackrel{}{x},t),^๐ช\stackrel{}{A}(\stackrel{}{x},t)=๐ช^T\stackrel{}{A}(๐ช\stackrel{}{x},t).$$
(7.4)
Here $`๐ช`$ is a general orthogonal $`3\times 3`$ rotation matrix with
$$๐ช^T\stackrel{}{\sigma }=g^{}\stackrel{}{\sigma }g,$$
(7.5)
where $`gSU(2)_s`$ represents the rotation $`๐ชSO(3)`$ in spinor space.
Frรถhlich and Studer noticed that the Pauli equation has a hidden local $`SU(2)_s`$ spin symmetry. This symmetry becomes manifest when one writes
$$iD_t\mathrm{\Psi }=\frac{1}{2M_e}D_iD_i\mathrm{\Psi },$$
(7.6)
with the $`SU(2)_sU(1)_Q`$ covariant derivative given by
$$D_\mu =_\mu +W_\mu (x)+ieA_\mu (x).$$
(7.7)
The components of the non-Abelian vector potential
$$W_\mu (x)=iW_\mu ^a(x)\frac{\sigma _a}{2},$$
(7.8)
can be identified as the electromagnetic field strengths $`\stackrel{}{E}(x)`$ and $`\stackrel{}{B}(x)`$, i.e.
$$W_t^a(x)=\mu _eB^a(x),W_i^a(x)=\frac{\mu _e}{2}\epsilon _{iab}E^b(x).$$
(7.9)
The anomalous magnetic moment $`\mu _e=g_ee/2M_e`$ of the electron (where, up to QED corrections, $`g_e=2`$) appears as a non-Abelian gauge coupling. The Abelian vector potential $`A_\mu (x)`$ is the usual one, except for a small contribution to the scalar potential due to the Darwin term,
$$A_t(x)=\mathrm{\Phi }(x)+\frac{1}{8M_e^2}\stackrel{}{}\stackrel{}{E}(x).$$
(7.10)
Hence, somewhat unexpected, the Pauli equation also transforms covariantly under local $`SU(2)_s`$ transformations
$$\mathrm{\Psi }(x)^{}=g(x)\mathrm{\Psi }(x),W_\mu (x)^{}=g(x)(W_\mu (x)+_\mu )g(x)^{}.$$
(7.11)
It should be pointed out that $`SU(2)_s`$ is not a gauge symmetry in the usual sense. In particular, the non-Abelian vector potential $`W_\mu (x)`$ is not an independent degree of freedom, but just given in terms of the external electromagnetic field strengths $`\stackrel{}{E}(x)`$ and $`\stackrel{}{B}(x)`$. The local $`SU(2)_s`$ symmetry is related to the global spatial rotations discussed before. In particular, global $`SU(2)_s`$ transformations take the form
$$\mathrm{\Psi }(x)^{}=g\mathrm{\Psi }(x),W_\mu (x)^{}=gW_\mu (x)g^{},$$
(7.12)
which, for example, implies
$$\stackrel{}{B}(x)^{}=๐ช^T\stackrel{}{B}(x),$$
(7.13)
where the resulting $`3\times 3`$ rotation matrix $`๐ชSO(3)`$ is again given by eq.(7.5). In contrast to a full spatial rotation, a global $`SU(2)_s`$ transformation does not rotate the argument $`\stackrel{}{x}`$ of the magnetic field to $`๐ช\stackrel{}{x}`$. Also the potentials $`\mathrm{\Phi }(x)`$ and $`\stackrel{}{A}(x)`$ are unaffected by the global $`SU(2)_s`$ symmetry. Consequently, the $`SU(2)_s`$ symmetry is inconsistent with the relations of eq.(7.2). Despite this, the local $`SU(2)_s`$ symmetry of the Pauli equation, which will be inherited by the Hubbard model and by the effective theory, dictates how low-frequency external electromagnetic fields are to be included in those theories. The high-frequency internal electromagnetic fields (for which eq.(7.2) is essential) are integrated out in the effective theory and thus do not spoil the symmetry. The local $`SU(2)_s`$ structure implies that in non-relativistic systems spin plays the role of an internal quantum number analogous to flavor in particle physics.
### 7.2 The Hubbard Model in an External Electromagnetic Field
In the next step we want to couple external electromagnetic fields to the Hubbard model. The Frรถhlich-Studer $`SU(2)_s`$ symmetry of the Pauli equation determines how to do this. One must simply use $`SU(2)_sU(1)_Q`$ covariant derivatives with $`\stackrel{}{E}(x)`$ and $`\stackrel{}{B}(x)`$ playing the role of non-Abelian vector potentials for $`SU(2)_s`$. Since the Hubbard model is defined on a spatial lattice, it is natural to construct corresponding $`SU(2)_sU(1)_Q`$ parallel transporters $`๐ฐ_{x,i}`$ connecting neighboring lattice sites $`x`$ and $`x+\widehat{i}`$,
$$๐ฐ_{x,i}=๐ซ\mathrm{exp}[_0^1๐sW_i(x+s\widehat{i})]\mathrm{exp}[ie_0^1๐sA_i(x+s\widehat{i})].$$
(7.14)
Here $`๐ซ`$ denotes path ordering along the link. Under local $`SU(2)_s`$ transformations the parallel transporter transforms as
$$๐ฐ_{x,i}^{}=g(x)๐ฐ_{x,i}g(x+\widehat{i})^{},$$
(7.15)
while under $`U(1)_Q`$ gauge transformations one has
$$^Q๐ฐ_{x,i}=\mathrm{exp}(i\omega (x))๐ฐ_{x,i}\mathrm{exp}(i\omega (x+\widehat{i})).$$
(7.16)
The Hubbard model Hamiltonian coupled to external electromagnetic fields then reads
$$H[๐ฐ]=t\underset{x,i}{}(c_x^{}๐ฐ_{x,i}c_{x+\widehat{i}}+c_{x+\widehat{i}}^{}๐ฐ_{x,i}^{}c_x)+\frac{U}{2}\underset{x}{}(c_x^{}c_x1)^2\mu \underset{x}{}(c_x^{}c_x1),$$
(7.17)
and the corresponding Schrรถdinger equation takes the form
$$iD_t\mathrm{\Psi }=H[๐ฐ]\mathrm{\Psi }.$$
(7.18)
Here $`\mathrm{\Psi }`$ is the multi-particle wave function and the covariant derivative is given by
$$D_t=_t+i\underset{x}{}[\stackrel{}{W}_t(x)\stackrel{}{S}_x+eA_t(x)Q_x].$$
(7.19)
It should be noted that the Zeeman coupling $`\mu _e\stackrel{}{B}(x)\stackrel{}{S}_x`$ enters the Hubbard model through $`D_t`$, while the spin-orbit coupling appears in the non-Abelian $`SU(2)_s`$ part of the parallel transporter $`๐ฐ_{x,i}`$.
In the Hilbert space of the theory local $`SU(2)_sU(1)_Q`$ transformations are implemented by unitary operators
$$V=\mathrm{exp}(i\underset{x}{}\stackrel{}{\eta }(x)\stackrel{}{S}_x),W=\mathrm{exp}(i\underset{x}{}\omega (x)Q_x),$$
(7.20)
such that
$`c_x^{}=V^{}c_xV=\mathrm{exp}(i\stackrel{}{\eta }(x){\displaystyle \frac{\stackrel{}{\sigma }}{2}})c_x=g(x)c_x,g(x)SU(2)_s,`$
$`{}_{}{}^{Q}c_{x}^{}=W^{}c_xW=\mathrm{exp}(i\omega (x))c_x,\mathrm{exp}(i\omega (x))U(1)_Q.`$ (7.21)
Together with eqs.(7.15) and (7.16) this implies that under the local transformations the Hamiltonian transforms as
$$H[๐ฐ^{}]=VH[๐ฐ]V^{},H[^Q๐ฐ]=WH[๐ฐ]W^{}.$$
(7.22)
Similarly, one obtains
$$D_t^{}=VD_tV^{},^QD_t=WD_tW^{},$$
(7.23)
such that the Schrรถdinger equation indeed transforms covariantly when one uses
$$\mathrm{\Psi }^{}=V\mathrm{\Psi },^Q\mathrm{\Psi }=W\mathrm{\Psi }.$$
(7.24)
### 7.3 External Electromagnetic Fields in the Effective Theory for Magnons and Charge Carriers
The couplings of magnons to external electromagnetic fields have been investigated in detail in . Again, the Frรถhlich-Studer symmetry is crucial and one obtains
$$S[\stackrel{}{e},W_\mu ]=d^2x๐t\frac{\rho _s}{2}\left(D_i\stackrel{}{e}D_i\stackrel{}{e}+\frac{1}{c^2}D_t\stackrel{}{e}D_t\stackrel{}{e}\right),$$
(7.25)
with the covariant derivative
$$D_\mu \stackrel{}{e}(x)=_\mu \stackrel{}{e}(x)+\stackrel{}{e}(x)\times \stackrel{}{W}_\mu (x).$$
(7.26)
Since magnons are electrically neutral, one may expect that they do not couple directly to the electromagnetic vector potential $`A_\mu (x)`$. Still, as discussed in the issue is potentially non-trivial because there is a Goldstone-Wilczek current
$$j_\mu ^{GW}(x)=\frac{1}{8\pi }\epsilon _{\mu \nu \rho }\stackrel{}{e}(x)[D_\nu \stackrel{}{e}(x)\times D_\rho \stackrel{}{e}(x)+\stackrel{}{W}_{\nu \rho }(x)],$$
(7.27)
with the non-Abelian field strength given by
$$\stackrel{}{W}_{\mu \nu }(x)=_\mu \stackrel{}{W}_\nu (x)_\nu \stackrel{}{W}_\mu (x)\stackrel{}{W}_\mu (x)\times \stackrel{}{W}_\nu (x).$$
(7.28)
The Goldstone-Wilczek current is an $`SU(2)_s`$ gauge-invariant extension of the baby-Skyrmion current of eq.(3.66) and is also topologically conserved, i.e. $`_\mu j_\mu ^{GW}=0`$. Hence, one may be tempted to add a Goldstone-Wilczek term $`j_\mu ^{GW}(x)A_\mu (x)`$ to the Lagrangian. However, just like the Hopf term, the Goldstone-Wilczek term breaks $`R`$, $`T`$, and $`T^{}`$ and is thus forbidden in the present case.
Using the $`P(x)`$ notation, in the presence of external electromagnetic fields the action of eq.(7.25) is given by
$$S[P,W_\mu ]=d^2x๐t\rho _s\left(\text{Tr}[D_iPD_iP]+\frac{1}{c^2}\text{Tr}[D_tPD_tP]\right),$$
(7.29)
where the $`SU(2)_s`$ covariant derivative is denoted by
$$D_\mu P(x)=_\mu P(x)+[W_\mu (x),P(x)].$$
(7.30)
As a consequence of the Frรถhlich-Studer symmetry, the action of eq.(7.29) is invariant even under local $`SU(2)_s`$ transformations
$$P(x)^{}=g(x)P(x)g(x)^{},W_\mu (x)^{}=g(x)(W_\mu (x)+_\mu )g(x)^{}.$$
(7.31)
Let us now discuss how external electromagnetic fields enter the fermionic part of the effective action. As a rule, ordinary derivatives must be replaced by covariant ones. This is the case also in the construction of the composite vector field which now takes the form
$$v_\mu (x)=u(x)D_\mu u(x)^{}=u(x)[_\mu +W_\mu (x)]u(x)^{}.$$
(7.32)
Under the local $`SU(2)_s`$ symmetry the field $`u(x)`$ transforms as
$$u(x)^{}=h(x)u(x)g(x)^{},$$
(7.33)
such that
$`v_\mu (x)^{}`$ $`=`$ $`h(x)u(x)g(x)^{}[_\mu +g(x)(W_\mu (x)+_\mu )g(x)^{}]g(x)u(x)^{}h(x)^{}`$ (7.34)
$`=`$ $`h(x)u(x)[_\mu +W_\mu (x)]u(x)^{}h(x)^{}`$
$`=`$ $`h(x)(v_\mu (x)+_\mu )h(x)^{}.`$
This is exactly the same transformation behavior as for the global $`SU(2)_s`$ transformation of eq.(3.42). In particular, this implies that the $`U(1)_s`$ covariant derivative $`D_\mu =_\mu +iv_\mu ^3(x)\sigma _3`$ need not be modified when $`SU(2)_s`$ is turned into a local symmetry. Of course, according to eq.(7.32), $`v_\mu (x)`$ now contains the electromagnetic fields $`\stackrel{}{E}(x)`$ and $`\stackrel{}{B}(x)`$ through the non-Abelian โgaugeโ field $`W_\mu (x)`$. Due to the local $`U(1)_Q`$ symmetry, the covariant derivatives still need to be extended to
$`D_\mu \mathrm{\Psi }^{A,B}=_\mu \mathrm{\Psi }^{A,B}+iv_\mu ^3(x)\sigma _3\mathrm{\Psi }^{A,B}+\mathrm{\Psi }^{A,B}ieA_\mu (x)\sigma _3,`$
$`D_\mu \mathrm{\Psi }^{A,B}=_\mu \mathrm{\Psi }^{A,B}\mathrm{\Psi }^{A,B}iv_\mu ^3(x)\sigma _3ieA_\mu (x)\sigma _3\mathrm{\Psi }^{A,B}.`$ (7.35)
It should also be noted that the low-energy effective theory is not necessarily just minimally coupled. In particular, the field strengths $`F_{\mu \nu }(x)=_\mu A_\nu (x)_\nu A_\mu (x)`$ and $`W_{\mu \nu }(x)`$ may also directly enter the low-energy effective theory.
## 8 Conclusions
We have constructed a systematic low-energy effective field theory describing the interactions of magnons with charge carriers doped into an antiferromagnet. A key ingredient for constructing the effective theory are symmetry considerations. The effective theory makes model-independent predictions for magnon-magnon, magnon-hole, and magnon-electron scattering. It also determines the long-range magnon-mediated forces between electrons or holes. Although these would be highly non-trivial non-perturbative issues from the point of view of Hubbard-type models, in the framework of the effective theory they can be understood quantitatively by perturbative analytic calculations. More ambitious non-perturbative questions might also be within reach of the effective theory. Such questions include the quantitative understanding of the Mott insulator state, the reduction of the staggered magnetization upon doping, the formation of a spiral phase, or the systematic investigation of dynamical mechanisms for the preformation of electron or hole pairs in the antiferromagnetic phase. In particular, magnon exchange โ the analog of pion exchange in nuclear physics โ suggests itself as a relevant mechanism.
Before one can do loop-calculations in the effective theory, one must establish a consistent power-counting scheme. This has originally been done for pion chiral perturbation theory , and carries over to magnon chiral perturbation theory in a straightforward manner. When charge carriers are included, the issue must be reconsidered. The same was true for baryon chiral perturbation theory of pions and nucleons. In the baryon number $`B=1`$ sector a consistent power-counting scheme enabling a systematic loop-expansion of the effective theory was established by Becher and Leutwyler . It is to be expected that this scheme can be extended to the low-energy theory of magnons and charge carriers developed here. The systematic power-counting in sectors with $`B2`$ still is a controversial issue in baryon chiral perturbation theory. The Weinberg power-counting scheme seems to work in most (but not necessarily in all) cases. Its relation to the alternative Kaplan-Savage-Wise scheme should be clarified further . In light of the experience with effective theories for the strong interactions, one should hence expect the issue of power-counting to be non-trivial in sectors with two or more charge carriers.
Even when the extra $`SU(2)_Q`$ symmetry is imposed, in the fermion sector the effective theory has a large number of low-energy parameters. There are two rest mass parameters $`M_1`$ and $`M_2`$ as well as two kinetic mass parameters $`M_1^{}`$ and $`M_2^{}`$ for the fermions, four coupling constants $`\mathrm{\Lambda }_1`$, $`\mathrm{\Lambda }_2`$, $`K_1`$, and $`K_2`$ for fermion-one-magnon vertices, two coupling constants $`N_1`$ and $`N_2`$ for fermion-two-magnon vertices, five 4-fermion coupling constants $`G_1,G_2,\mathrm{},G_5`$, two 6-fermion couplings $`H_1`$ and $`H_2`$, and finally one 8-fermion coupling $`I`$. If only the $`U(1)_Q`$ symmetry is imposed there are even more parameters. The large number of a priori undetermined low-energy parameters is the price one has to pay for the universality and model-independence of the effective theory. Only in this way the low-energy physics of any arbitrary cuprate antiferromagnet can be captured by the effective theory. Of course, due to the rather large number of parameters, the predictive power of the effective theory is somewhat limited. Still, only a few parameters enter in some relevant physical quantities. For example, the one-magnon exchange potential between charge carriers depends only on certain combinations of the fermion-magnon couplings $`\mathrm{\Lambda }_1`$, $`\mathrm{\Lambda }_2`$, $`K_1`$, and $`K_2`$. Also, for example, the details of the short-range 4-, 6-, and 8-fermion couplings are not expected to be essential for identifying potential mechanisms for preforming electron or hole pairs in the antiferromagnetic phase. It is interesting to note that the low-energy effective theory of the $`t`$-$`J`$ model, in which electrons are excluded beyond half-filling and holes are the only charge carriers, has a much smaller number of low-energy parameters. In that case, there are only one rest mass parameter $`M`$, one kinetic mass parameter $`M^{}`$, two coupling constants $`\mathrm{\Lambda }`$ and $`K`$ for hole-one-magnon vertices, one coupling constant $`N`$ for a hole-two-magnon vertex, and one 4-fermion coupling constant $`G`$. It would be interesting to perform numerical simulations of the Hubbard or $`t`$-$`J`$ model in order to determine the values of the corresponding low-energy parameters by comparison with calculations in the effective theory. For example, in the $`t`$-$`J`$ model one can determine the parameters $`M`$, $`M^{}`$, $`\mathrm{\Lambda }`$, $`K`$, and $`N`$ from simulations in the one-hole sector, while the determination of $`G`$ requires computations in the two-hole sector of the Hilbert space.
It should be pointed out that, as it stands, the effective theory is applicable only at small doping, i.e. for small $`\mu `$. This is sufficient for understanding the long-range forces between electrons or holes in the antiferromagnetic phase. It should also allow a quantitative investigation of the reduction of the staggered magnetization upon doping. However, in order to enter the high-temperature superconducting phase itself, if this is at all possible within the effective field theory presented here, larger values of $`\mu `$ will be necessary. Once $`\mu `$ becomes large, it sets a new scale which must be taken into account in the power-counting. However, most important, the symmetry considerations of the present paper still apply in that case as well.
Some of the most interesting questions one can address in the framework of the effective theory may require non-perturbative calculations. While in some cases such calculations can be performed in the continuum, in others they may require a non-perturbative regularization of the effective theory. In the effective theory of pions and nucleons was regularized on a space-time lattice in order to address non-perturbative questions concerning the strong interactions. It may also be useful to formulate the effective theory of magnons and charge carriers on the lattice. For example, it would be interesting to investigate if the effective theory is more easily solvable by numerical simulation than the standard Hubbard-type models.
We like to emphasize again that effective field theory also allows us to include phonons in addition to magnons. This may shed light on more complicated potential mechanisms for Cooper pair preformation which involve both magnon and phonon exchange. It is interesting to construct such an effective theory. In particular, the Galilean (or even Poincarรฉ) symmetry is then non-linearly realized.
To summarize, low-energy effective field theory is a powerful tool that has several advantages compared to the direct use of microscopic models. First, it is model-independent and provides universal predictions. Material-specific details of the underlying microscopic system enter the effective theory only through low-energy parameters whose values can be determined by comparison with experiments or with numerical simulations. Second, and most important, the electrons or holes of the effective theory are quasi-particles whose long-range forces are weak and calculable in perturbation theory. This is a significant advantage compared to calculations in microscopic models of strongly correlated electrons which are necessarily non-perturbative. While it is practically impossible to reliably determine the long-range forces between charge carriers from Hubbard-type models, in the effective theory the calculation of the one-magnon exchange forces is straightforward and presently in progress. It is very interesting to ask if these forces will provide a potential mechanism for the preformation of electron or hole pairs. In any case, we propose the systematic low-energy effective field theory approach as a better compromise between calculability and predictive power than the one offered by Hubbard-type models. Effective field theory sheds new light on the dynamics of charge carriers in antiferromagnets, and there is hope that it may even be applicable to the high-temperature superconductors themselves.
## Acknowledgements
We have benefitted from discussions with M. Bissegger, S. Chandrasekharan, G. Colangelo, J. Gasser, P. Hasenfratz, H. Leutwyler, P. Minkowski, and F. Niedermayer. This work is supported by funds provided by the Schweizerischer Nationalfonds.
## Appendix A Electron-Hole Representation of the Hubbard Model Operators
For $`U|t|`$ the Hubbard model at half-filling reduces to the antiferromagnetic quantum Heisenberg model. In contrast to the Heisenberg ferromagnet, the ground state of the antiferromagnet is not known analytically. In particular, the naive Nรฉel state
$$|N=\underset{xA}{}c_x^{}\underset{xB}{}c_x^{}|0,$$
(A.1)
with all spins down on the even sublattice $`A`$ and all spins up on the odd sublattice $`B`$ is not an eigenstate of the Hubbard Hamiltonian. Still, we use this state in order to define electron and hole operators. For even sites we then find
$$c_x|N=0,c_x^{}|N=0,xA.$$
(A.2)
Correspondingly, $`c_x^{}`$ creates an electron, while $`c_x`$ creates a hole. Hence, just like a relativistic Dirac spinor, the $`SU(2)_s`$ spinor
$$c_x=\left(\begin{array}{c}c_x\\ c_x\end{array}\right)=\left(\begin{array}{c}a_x\\ b_x^{}\end{array}\right),xA,$$
(A.3)
consists of a particle annihilation operator $`a_x`$ in the upper component and a hole creation operator $`b_x^{}`$ in the lower component. Note that the annihilation of an electron with spin down via $`c_x`$ corresponds to the creation of a hole with spin up via $`b_x^{}`$. Similarly, on the odd sites one has
$$c_x|N=0,c_x^{}|N=0,xB.$$
(A.4)
In this case, $`c_x^{}`$ creates a particle, while $`c_x`$ creates a hole and we write
$$c_x=\left(\begin{array}{c}c_x\\ c_x\end{array}\right)=\left(\begin{array}{c}b_x^{}\\ a_x\end{array}\right),xB.$$
(A.5)
## Appendix B Removal of Non-Canonical Terms by a Field Redefinition
The most general $`SU(2)_Q`$-breaking but $`U(1)_Q`$-symmetric terms containing one covariant time-derivative are given by
$`{\displaystyle \frac{a}{2}}\text{Tr}[\mathrm{\Psi }^AD_t\mathrm{\Psi }^A+\mathrm{\Psi }^BD_t\mathrm{\Psi }^B]+{\displaystyle \frac{b}{2}}\text{Tr}[\mathrm{\Psi }^A\sigma _3D_t\mathrm{\Psi }^A\sigma _3\mathrm{\Psi }^B\sigma _3D_t\mathrm{\Psi }^B\sigma _3]`$
$`+{\displaystyle \frac{c}{2}}\text{Tr}[\mathrm{\Psi }^AD_t\mathrm{\Psi }^B\sigma _3+\mathrm{\Psi }^BD_t\mathrm{\Psi }^A\sigma _3]`$
$`=(\psi _+^A,\psi _+^B)\left(\begin{array}{cc}a+b& c\\ c& ab\end{array}\right)\left(\begin{array}{c}D_t\psi _+^A\\ D_t\psi _+^B\end{array}\right)`$ (B.5)
$`+(\psi _{}^A,\psi _{}^B)\left(\begin{array}{cc}ab& c\\ c& a+b\end{array}\right)\left(\begin{array}{c}D_t\psi _{}^A\\ D_t\psi _{}^B\end{array}\right)`$ (B.10)
$`=(\stackrel{~}{\psi }_+^A,\stackrel{~}{\psi }_+^B)\left(\begin{array}{c}D_t\stackrel{~}{\psi }_+^A\\ D_t\stackrel{~}{\psi }_+^B\end{array}\right)+(\stackrel{~}{\psi }_{}^A,\stackrel{~}{\psi }_{}^B)\left(\begin{array}{c}D_t\stackrel{~}{\psi }_{}^A\\ D_t\stackrel{~}{\psi }_{}^B\end{array}\right).`$ (B.15)
Here $`\stackrel{~}{\psi }_\pm ^{A,B}(x)`$ results from a field redefinition that diagonalizes the matrices in the previous expression. Only the term proportional to $`a`$ contains the standard form $`\psi _+^A_t\psi _+^A+\psi _{}^A_t\psi _{}^A+\psi _+^B_t\psi _+^B+\psi _{}^B_t\psi _{}^B`$ which implies canonical anticommutation relations between fermionic creation and annihilation operators in the Hamiltonian formulation. The non-canonical terms (proportional to $`b`$ and $`c`$) can be removed by an appropriate field redefinition
$`\left(\begin{array}{c}\stackrel{~}{\psi }_\pm ^A(x)\\ \stackrel{~}{\psi }_\pm ^B(x)\end{array}\right)=\left(\begin{array}{cc}\sqrt{\lambda _\pm }& 0\\ 0& \sqrt{\lambda _{}}\end{array}\right)U_\pm \left(\begin{array}{c}\psi _\pm ^A(x)\\ \psi _\pm ^B(x)\end{array}\right),`$ (B.22)
$`\lambda _\pm =a\pm \sqrt{b^2+c^2},U_\pm =\left(\begin{array}{cc}X& \pm Y\\ Y& X\end{array}\right).`$ (B.25)
Here $`U_\pm `$ are unitary matrices with $`X,Y๐จ๐ฑ`$ which obey
$$U_\pm \left(\begin{array}{cc}a\pm b& c\\ c& ab\end{array}\right)U_\pm ^{}=\left(\begin{array}{cc}\lambda _\pm & 0\\ 0& \lambda _{}\end{array}\right).$$
(B.26)
It is straightforward to show that the redefined fields $`\stackrel{~}{\psi }_\pm ^{A,B}(x)`$ have the same symmetry properties of eqs.(5.1) and (5.1) as the original fields $`\psi _\pm ^{A,B}(x)`$. Under the $`SU(2)_s`$ symmetry the original fields transform as
$$\psi _\pm ^{A,B}(x)^{}=\mathrm{exp}(\pm i\alpha (x))\psi _\pm ^{A,B}(x),$$
(B.27)
and after the field redefinition again
$`\stackrel{~}{\psi }_\pm ^A(x)^{}`$ $`=`$ $`\sqrt{\lambda _\pm }[X\psi _\pm ^A(x)^{}\pm Y\psi _\pm ^B(x)^{}]`$
$`=`$ $`\mathrm{exp}(\pm i\alpha (x))\sqrt{\lambda _\pm }[X\psi _\pm ^A(x)\pm Y\psi _\pm ^B(x)]=\mathrm{exp}(\pm i\alpha (x))\stackrel{~}{\psi }_\pm ^A(x),`$
$`\stackrel{~}{\psi }_\pm ^B(x)^{}`$ $`=`$ $`\sqrt{\lambda _{}}[Y\psi _\pm ^A(x)^{}+X\psi _\pm ^B(x)^{}]`$ (B.28)
$`=`$ $`\mathrm{exp}(\pm i\alpha (x))\sqrt{\lambda _{}}[Y\psi _\pm ^A(x)+X\psi _\pm ^B(x)]=\mathrm{exp}(\pm i\alpha (x))\stackrel{~}{\psi }_\pm ^B(x).`$
Similarly, under the $`U(1)_Q`$ symmetry the original fields transform as
$${}_{}{}^{Q}\psi _{\pm }^{A,B}(x)=\mathrm{exp}(i\omega )\psi _\pm ^{A,B}(x),$$
(B.29)
and again
$`{}_{}{}^{Q}\stackrel{~}{\psi }_{\pm }^{A}(x)`$ $`=`$ $`\sqrt{\lambda _\pm }[X^Q\psi _\pm ^A(x)\pm Y^Q\psi _\pm ^B(x)]`$
$`=`$ $`\mathrm{exp}(i\omega )\sqrt{\lambda _\pm }[X\psi _\pm ^A(x)\pm Y\psi _\pm ^B(x)]=\mathrm{exp}(i\omega )\stackrel{~}{\psi }_\pm ^A(x),`$
$`{}_{}{}^{Q}\stackrel{~}{\psi }_{\pm }^{B}(x)`$ $`=`$ $`\sqrt{\lambda _{}}[Y^Q\psi _\pm ^A(x)+X^Q\psi _\pm ^B(x)]`$ (B.30)
$`=`$ $`\mathrm{exp}(i\omega )\sqrt{\lambda _{}}[Y\psi _\pm ^A(x)+X\psi _\pm ^B(x)]=\mathrm{exp}(i\omega )\stackrel{~}{\psi }_\pm ^B(x).`$
Under the modified displacement symmetry $`D^{}`$ one has
$${}_{}{}^{D^{}}\psi _{\pm }^{A,B}(x)=\pm \psi _{}^{B,A}(x),$$
(B.31)
and after the field redefinition one again obtains
$`{}_{}{}^{D^{}}\stackrel{~}{\psi }_{\pm }^{A}(x)`$ $`=`$ $`\sqrt{\lambda _\pm }[X^D^{}\psi _\pm ^A(x)\pm Y^D^{}\psi _\pm ^B(x)]`$
$`=`$ $`\pm \sqrt{\lambda _\pm }[X\psi _{}^B(x)\pm Y\psi _{}^A(x)]=\pm \stackrel{~}{\psi }_{}^B(x),`$
$`{}_{}{}^{D^{}}\stackrel{~}{\psi }_{\pm }^{B}(x)`$ $`=`$ $`\sqrt{\lambda _{}}[Y^D^{}\psi _\pm ^A(x)+X^D^{}\psi _\pm ^B(x)]`$ (B.32)
$`=`$ $`\pm \sqrt{\lambda _{}}[Y\psi _{}^B(x)+X\psi _{}^A(x)]=\pm \stackrel{~}{\psi }_{}^A(x).`$
Since the displacement symmetry $`D`$ is a combination of $`D^{}`$ and $`SU(2)_s`$ it also maintains its original form. The same is true for the discrete symmetries $`O`$ and $`R`$. Finally, under the modified time-reversal $`T^{}`$ the original fields transform as
$${}_{}{}^{T^{}}\psi _{\pm }^{A,B}(x)=\psi _\pm ^{A,B}(Tx),$$
(B.33)
such that
$`{}_{}{}^{T^{}}\stackrel{~}{\psi }_{\pm }^{A}(x)`$ $`=`$ $`\sqrt{\lambda _\pm }[X^T^{}\psi _\pm ^A(x)\pm Y^T^{}\psi _\pm ^B(x)]`$
$`=`$ $`\sqrt{\lambda _\pm }[X\psi _\pm ^A(Tx)\pm Y\psi _\pm ^B(Tx)]=\stackrel{~}{\psi }_\pm ^A(Tx),`$
$`{}_{}{}^{T^{}}\stackrel{~}{\psi }_{\pm }^{B}(x)`$ $`=`$ $`\sqrt{\lambda _{}}[Y^T^{}\psi _\pm ^A(x)+X^T^{}\psi _\pm ^B(x)]`$ (B.34)
$`=`$ $`\sqrt{\lambda _{}}[Y\psi _\pm ^A(Tx)+X\psi _\pm ^B(Tx)]=\stackrel{~}{\psi }_\pm ^B(Tx).`$
As a combination of $`T^{}`$ and $`SU(2)_s`$ the time-reversal symmetry $`T`$ also maintains its original form after the field redefinition. The only symmetry that does not maintain its original form is $`SU(2)_Q`$. This is no problem since the non-canonical terms can arise only when the $`SU(2)_Q`$ symmetry is explicitly broken down to $`U(1)_Q`$ and is hence no longer a symmetry of the theory.
Since the redefined fields transform exactly like the original ones, the terms in the effective Lagrangian take exactly the same form as before. Hence, it is indeed justified not to include the non-canonical terms in the effective Lagrangian.
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# New Classes of Almost Bent and Almost Perfect Nonlinear Polynomials11footnote 1Part of this paper was presented at WCC 2005 [4].
## 1 Introduction
Vectorial Boolean functions are used in cryptography, more precisely in block ciphers. An important condition on these functions is a high resistance to the differential and linear cryptanalyses, which are the main attacks on block ciphers. The functions with the smallest possible differential uniformity (see ) oppose an optimum resistance to the differential attack . They are called almost perfect nonlinear (APN). The functions achieving the maximal possible nonlinearity (see ) possess the best resistance to the linear attack and they are called almost bent (AB) or maximum nonlinear.
Up to now, in the study of APN and AB functions the main attention has been payed to power mappings and all known constructions of APN and AB functions happen to be equivalent to power functions (see for example ). Recall that functions $`F`$ and $`F^{}`$ are called equivalent if either $`F^{}=A_1FA_2+A`$ or $`F^{}=A_1F^1A_2+A`$ (in case $`F`$ is a permutation) for some affine functions $`A`$, $`A_1`$ and $`A_2`$, where $`A_1`$ and $`A_2`$ are permutations. We shall say that two functions $`F`$ and $`F^{}=A_1FA_2+A`$ (with $`A,A_1,A_2`$ affine, $`A_1,A_2`$ being permutations) are extended affine equivalent (EA-equivalent). In this paper we give the first theoretical constructions of APN and AB polynomials which are inequivalent to power mappings. In these constructions, we apply the transformation of functions given by Carlet, Charpin and Zinoviev (see Proposition 3 in ) to the Gold APN and AB mappings . This transformation leads to an equivalence relation of functions that we call the Carlet-Charpin-Zinoviev equivalence (CCZ-equivalence). CCZ-equivalence corresponds to the affine equivalence of the graphs of functions, i.e. functions $`F`$ and $`F^{}`$ are CCZ-equivalent if and only if, for some affine permutation, the image of the graph of $`F`$ is the graph of the function $`F^{}`$. CCZ-equivalence preserves the nonlinearity, the differential uniformity of functions and the resistance of a function to the algebraic cryptanalysis . It could be expected that CCZ-equivalence coincides in practice with the equivalence mentioned above. The present paper aims at proving and illustrating that CCZ-equivalence is more general.
The next section contains all the necessary definitions related to vectorial Boolean functions, including EA-equivalence, APN and AB properties.
In Section 3, we give the definition of CCZ-equivalence, we describe its main properties and we show its connections with EA-equivalence.
We give some results related to a classification of functions CCZ-equivalent to the Gold mappings in Section 4.
Theorems 1, 2 and 4 in Section 5 present constructions of AB and APN polynomials which are EA-inequivalent to power functions and Theorem 3 presents an APN function EA-inequivalent to all known APN mappings.
This paper is an extended version of an abstract published in . In particular, Theorem 4, the proof of Theorem 3 and some properties of CCZ-equivalence are not included in the abstract.
## 2 Almost perfect nonlinear and almost bent functions
Let $`๐ฝ_2^m`$ be the $`m`$-dimensional vector space over the field $`๐ฝ_2`$. Any function $`F`$ from $`๐ฝ_2^m`$ to itself can be uniquely represented as a polynomial on $`m`$ variables with coefficients in $`๐ฝ_2^m`$, whose degree with respect to each coordinate is at most 1:
$$F(x_1,\mathrm{},x_m)=\underset{u๐ฝ_2^m}{}c(u)\left(\underset{i=1}{\overset{m}{}}x_i^{u_i}\right),c(u)๐ฝ_2^m.$$
This representation is called the *algebraic normal form* of $`F`$ and its degree $`d^{}(F)`$ the *algebraic degree* of the function $`F`$.
Besides, viewed as a function from the field $`๐ฝ_{2^m}`$ to itself, $`F`$ has a unique representation as a univariate polynomial over $`๐ฝ_{2^m}`$ of degree smaller than $`2^m`$:
$$F(x)=\underset{i=0}{\overset{2^m1}{}}c_ix^i,c_i๐ฝ_{2^m}.$$
For any $`k`$, $`0k2^m1`$, the number $`w_2(k)`$ of the nonzero coefficients $`k_s\{0,1\}`$ in the binary expansion $`_{s=0}^{m1}2^sk_s`$ of $`k`$ is called the $`2`$-weight of $`k`$. The algebraic degree of $`F`$ is equal to the maximum 2-weight of the exponents $`i`$ of the polynomial $`F(x)`$ such that $`c_i0`$, that is $`d^{}(F)=\mathrm{max}_{0in1,c_i0}w_2(i)`$ (see ).
A function $`F:๐ฝ_2^m๐ฝ_2^m`$ is *linear* if and only if $`F(x)`$ is a linearized polynomial over $`๐ฝ_{2^m}`$, that is,
$$\underset{i=0}{\overset{m1}{}}c_ix^{2^i},c_i๐ฝ_{2^m}.$$
The sum of a linear function and a constant is called an *affine function*.
Let $`F`$ be a function from $`๐ฝ_2^m`$ to itself and $`A_1`$, $`A_2:๐ฝ_2^m๐ฝ_2^m`$ be affine permutations. Then the functions $`F`$ and $`A_1FA_2`$ are called *affine equivalent*. Affine equivalent functions have the same algebraic degree (i.e. the algebraic degree is *affine invariant*).
We shall say that the functions $`F`$ and $`F^{}`$ are *extended affine equivalent* (EA-equivalent) if $`F^{}=A_1FA_2+A`$ for some affine permutations $`A_1`$, $`A_2`$ and an affine function $`A`$. If $`F`$ is not affine, then $`F`$ and $`F^{}`$ have again the same algebraic degree.
For a function $`F:๐ฝ_2^m๐ฝ_2^m`$ and any elements $`a,b๐ฝ_2^m`$ we denote
$$\delta _F(a,b)=|\{x๐ฝ_2^m:F(x+a)+F(x)=b\}|$$
and
$$\mathrm{\Delta }_F=\{\delta _F(a,b):a,b๐ฝ_2^m,a0\}.$$
$`F`$ is called a *differentially $`\delta `$-uniform* function if $`\mathrm{max}_{a๐ฝ_2^m,b๐ฝ_2^m}\delta _F(a,b)\delta ,`$ where $`๐ฝ_2^m=๐ฝ_2^m\{0\}`$. For any $`a,b๐ฝ_2^m`$, the number $`\delta _F(a,b)`$ is even since if $`x_0`$ is a solution of the equation $`F(x+a)+F(x)=b`$ then $`x_0+a`$ is a solution too. Hence, $`\delta 2`$. Differentially 2-uniform mappings are called *almost perfect nonlinear*.
For any function $`F:๐ฝ_2^m๐ฝ_2^m`$, the values
$$\lambda _F(a,b)=\underset{x๐ฝ_2^m}{}(1)^{bF(x)+ax},a,b๐ฝ_2^m,$$
do not depend on a particular choice of the inner product $`\mathrm{"}\mathrm{"}`$ in $`๐ฝ_2^m`$. If we identify $`๐ฝ_2^m`$ with $`๐ฝ_{2^m}`$ then we can take $`xy=tr(xy)`$, where $`tr(x)=x+x^2+x^4+\mathrm{}+x^{2^{m1}}`$ is the trace function from $`๐ฝ_{2^m}`$ into $`๐ฝ_2`$. Later we shall also need for a divisor $`n`$ of $`m`$ the relative trace $`tr_{m/n}(x)=x+x^{2^n}+x^{2^{2n}}+\mathrm{}+x^{2^{(m/n1)n}}`$ and we denote $`tr_n(x)=x+x^2+x^{2^2}+\mathrm{}+x^{2^{n1}}`$. The set $`\mathrm{\Lambda }_F=\{\lambda _F(a,b):a,b๐ฝ_2^m,b0\}`$ is called the *Walsh spectrum* of $`F`$ and the value
$$๐ฉ(F)=2^{m1}\frac{1}{2}\underset{a๐ฝ_2^m,b๐ฝ_2^m}{\mathrm{max}}|\lambda _F(a,b)|$$
the *nonlinearity* of the function $`F`$. The nonlinearity of any function $`F`$ satisfies the inequality
$$๐ฉ(F)2^{m1}2^{\frac{m1}{2}}$$
() and in case of equality $`F`$ is called *almost bent* or *maximum nonlinear*. For any AB function $`F`$, the Walsh spectrum $`\mathrm{\Lambda }_F`$ equals $`\{0,\pm 2^{\frac{m+1}{2}}\}`$ as it is proven in .
For EA-equivalent functions $`F`$ and $`F^{}`$, we have $`\mathrm{\Delta }_F=\mathrm{\Delta }_F^{}`$, $`\mathrm{\Lambda }_F=\mathrm{\Lambda }_F^{}`$ and if $`F`$ is a permutation then $`\mathrm{\Delta }_F=\mathrm{\Delta }_{F^1}`$, $`\mathrm{\Lambda }_F=\mathrm{\Lambda }_{F^1}`$ (see ). Therefore, if $`F`$ is APN (resp. AB) and $`F^{}`$ is EA-equivalent to either $`F`$ or $`F^1`$ (if $`F`$ is a permutation), then $`F^{}`$ is also APN (resp. AB).
Table 1 (resp. Table 2) gives all known values of exponents $`d`$ (up to EA-equivalence and up to taking the inverse when a function is a permutation) such that the power function $`x^d`$ is APN (resp. AB).
Table 1
Known APN power functions $`x^d`$ on $`๐ฝ_{2^m}`$.
| | Exponents $`d`$ | Conditions | $`w_2\left(d\right)`$ | Proven in |
| --- | --- | --- | --- | --- |
| Gold functions | $`2^i+1`$ | $`\mathrm{gcd}(i,m)=1`$ | 2 | |
| Kasami functions | $`2^{2i}2^i+1`$ | $`\mathrm{gcd}(i,m)=1`$ | $`i+1`$ | |
| Welch function | $`2^t+3`$ | $`m=2t+1`$ | 3 | |
| Niho function | $`2^t+2^{\frac{t}{2}}1`$, $`t`$ even | $`m=2t+1`$ | $`\left(t+2\right)/2`$ | |
| | $`2^t+2^{\frac{3t+1}{2}}1`$, $`t`$ odd | | $`t+1`$ | |
| Inverse function | $`2^{2t}1`$ | $`m=2t+1`$ | $`m1`$ | |
| Dobbertin function | $`2^{4i}+2^{3i}+2^{2i}+2^i1`$ | $`m=5i`$ | $`i+3`$ | |
Table 2
Known AB power functions $`x^d`$ on $`๐ฝ_{2^m}`$, $`m`$ odd.
| | Exponents $`d`$ | Conditions | Proven in |
| --- | --- | --- | --- |
| Gold functions | $`2^i+1`$ | $`\mathrm{gcd}(i,m)=1`$ | |
| Kasami functions | $`2^{2i}2^i+1`$ | $`\mathrm{gcd}(i,m)=1`$ | |
| Welch function | $`2^t+3`$ | $`m=2t+1`$ | |
| Niho function | $`2^t+2^{\frac{t}{2}}1`$, $`t`$ even | $`m=2t+1`$ | |
| | $`2^t+2^{\frac{3t+1}{2}}1`$, $`t`$ odd | | |
Every AB function is APN . The converse is not true since AB functions exist only when $`m`$ is odd while APN functions exist for $`m`$ even too. Besides, in the $`m`$ odd case, the Dobbertin APN function is not AB as proven in . Also the inverse APN function is not AB since it has the algebraic degree $`m1`$ while the algebraic degree of any AB function is not greater than $`(m+1)/2`$ (see ). When $`m`$ is even, the inverse function $`x^{2^m2}`$ is a differentially 4-uniform permutation and has the best known nonlinearity , that is $`2^{m1}2^{\frac{m}{2}}`$ (see ). When $`m2`$ \[mod 4\] and $`\mathrm{gcd}(i,m)=2`$, the functions $`x^{2^i+1}`$ and $`x^{2^{2i}2^i+1}`$ also have the best known nonlinearity and when $`\mathrm{gcd}(i,m)=s`$ and $`m/s`$ is odd, these functions have three valued Walsh spectrum $`\{0,\pm 2^{\frac{m+s}{2}}\}`$ ().
## 3 Carlet-Charpin-Zinoviev equivalence of functions
For a function $`F`$ from $`๐ฝ_2^m`$ to itself, we denote by $`G_F`$ the *graph of the function* $`F`$:
$$G_F=\{(x,F(x)):x๐ฝ_2^m\}๐ฝ_2^{2m}.$$
The property of stability of APN and AB functions given in Proposition 3 of leads to the definition of the following equivalence relation of functions.
###### Definition 1
*(CCZ-equivalence)* We say that functions $`F,F^{}:๐ฝ_2^m๐ฝ_2^m`$ are *Carlet-Charpin-Zinoviev equivalent* if there exists a linear (affine) permutation $`:๐ฝ_2^{2m}๐ฝ_2^{2m}`$ such that $`(G_F)=G_F^{}`$.
We shall consider only the case of linear functions, but all proofs for the statements related to the CCZ-equivalence can be easily extended for the case of affine functions too.
A linear function $`:๐ฝ_2^{2m}๐ฝ_2^{2m}`$ can be considered as a pair of linear functions $`L_1,L_2:๐ฝ_2^{2m}๐ฝ_2^m`$ such that $`(x,y)=(L_1(x,y),L_2(x,y))`$ for all $`x,y๐ฝ_2^m`$. Then for a function $`F:๐ฝ_2^m๐ฝ_2^m`$ we have $`(x,F(x))=(F_1(x),F_2(x))`$, where
$$F_1(x)=L_1(x,F(x)),$$
(1)
$$F_2(x)=L_2(x,F(x)).$$
(2)
Hence, the set $`(G_F)=\{(F_1(x),F_2(x)):x๐ฝ_2^m\}`$ is the graph of a function $`F^{}`$ if and only if the function $`F_1`$ is a permutation. If $``$ and $`F_1`$ are permutations then $`(G_F)=G_F^{}`$, where $`F^{}=F_2F_1^1`$, and the functions $`F`$ and $`F^{}`$ are CCZ-equivalent.
In the proposition below, we give a slightly different approach to the CCZ-equivalence, that we shall use in the constructions of APN polynomials which will be EA-inequivalent to power functions. Recall that a set $`G๐ฝ_2^{2m}`$ is transversal to a subgroup $`V`$ of $`๐ฝ_2^{2m}`$ if $`|G(u+V)|=1`$ for any $`u๐ฝ_2^{2m}`$.
###### Proposition 1
Let $`F:๐ฝ_2^m๐ฝ_2^m`$ and $`G๐ฝ_2^{2m}`$. Then the set $`G`$ is the graph of a function CCZ-equivalent to $`F`$ if and only if there exists a linear permutation $`:๐ฝ_2^{2m}๐ฝ_2^{2m}`$ such that $`G=(G_F)`$ and $`G_F`$ is transversal to $`V^{}=^1(V)`$, where $`V=\{(0,x):x๐ฝ_2^m\}.`$
*Proof.* The condition that there exists a linear permutation $`:๐ฝ_2^{2m}๐ฝ_2^{2m}`$ such that $`G=(G_F)`$ is clearly necessary. Let us denote $`U=\{(x,0):x๐ฝ_2^m\}`$, $`V=\{(0,x):x๐ฝ_2^m\}`$, $`^1(U)=U^{}`$ and $`^1(V)=V^{}`$. Then $`G`$ is the graph of a function if and only if $`|G(u+V)|=1`$ for any $`uU`$; that is, if and only if $`|^1(G)(u^{}+V^{})|=1`$ for any $`u^{}U^{}`$. Hence, $`G`$ is the graph of a function CCZ-equivalent to $`F`$ if and only if $`G_F`$ is transversal to $`V^{}`$. $`\mathrm{}`$
Let $`F`$ be an arbitrary function on $`๐ฝ_2^m`$ and $`V^{}`$ be an arbitrary subgroup of $`๐ฝ_2^{2m}`$. If we denote
$$H_a=\{F(x+a)+F(x):x๐ฝ_2^m\},a๐ฝ_2^m,$$
$$A_F=\{(a,F(x+a)+F(x)):a๐ฝ_{2^m}^{},x๐ฝ_{2^m}\},$$
then $`A_F=_{a๐ฝ_{2^m}^{}}(a,H_a).`$ It is easy to note that $`A_FV^{}=\mathrm{}`$ if and only if $`G_F`$ is transversal to $`V^{}`$. We shall see that, to construct a function CCZ-equivalent but EA-inequivalent to the function $`F`$, it is necessary that the subgroup $`V^{}`$ in Proposition 1 is different from $`V`$ and as soon as we have a subroup $`V^{}`$ such that $`G_F`$ is transversal to $`V^{}`$ we can find a linear permutation $``$ that $`V^{}=^1(V)`$. However, even if we have such a subgroup it is difficult to determine whether the resulting function is EA-inequivalent to $`F`$.
The property of stability of APN and AB mappings (see ) can be easily generalized to all functions (not necessarily APN or AB) as follows:
###### Proposition 2
Let $`F,F^{}:๐ฝ_2^m๐ฝ_2^m`$ be CCZ-equivalent functions. Then $`\mathrm{\Delta }_F=\mathrm{\Delta }_F^{}`$ and $`\mathrm{\Lambda }_F=\mathrm{\Lambda }_F^{}`$. In particular, $`F`$ is an APN (resp. AB) function if and only if $`F^{}`$ is APN (resp. AB).
*Proof.* If $`F,F^{}:๐ฝ_2^m๐ฝ_2^m`$ are CCZ-equivalent, then $`F^{}=F_2F_1^1`$ for a certain linear permutation $`=(L_1,L_2)`$, where $`F_1,F_2`$ are defined by (1) and (2). For any $`(a,b)๐ฝ_2^{2m}`$ we have
$`\lambda _F(a,b)`$ $`=`$ $`{\displaystyle \underset{x๐ฝ_2^m}{}}(1)^{bF(x)+ax}={\displaystyle \underset{x๐ฝ_2^m}{}}(1)^{(a,b)(x,F(x))}`$
$`=`$ $`{\displaystyle \underset{x๐ฝ_2^m}{}}(1)^{(a,b)^1(F_1(x),F_2(x))}={\displaystyle \underset{x๐ฝ_2^m}{}}(1)^{^1(a,b)(x,F_2F_1^1(x))}`$
$`=`$ $`\lambda _F^{}(^1(a,b)),`$
where $`^1`$ is the adjoint operator of $`^1`$ (i.e. $`x^1(y)=^1(x)y`$, for any $`(x,y)๐ฝ_2^{2m}`$ ; if โ$``$โ is the usual inner product, then $`^1`$ is the linear permutation whose matrix is transposed of that of $`^1`$). Hence, $`\mathrm{\Lambda }_F=\mathrm{\Lambda }_F^{}`$.
The proof that $`\mathrm{\Delta }_F=\mathrm{\Delta }_F^{}`$ for arbitrary functions $`F,F^{}`$ is completely the same as in the case when one of the functions $`F,F^{}`$ is APN (see ). $`\mathrm{}`$
Remark 1 Obviously, CCZ-equivalence can be defined for functions between any two groups $`H_1`$ and $`H_2`$. For a function $`F:H_1H_2`$ we can consider the set of the values $`\delta _F(a,b)=|\{xH_1:F(x+a)F(x)=b\}|`$, $`aH_1\backslash \{0\}`$, $`bH_2`$, and if the groups $`H_1`$ and $`H_2`$ are Abelian, then the discrete Fourier transform of $`F`$ can also be defined. In this more general case CCZ-equivalent functions again have the same differential and linear properties. One can find results related to nonlinear functions on Abelian groups in . $`\mathrm{}`$
Since CCZ-equivalent functions have the same differential uniformity and the same nonlinearity, the resistance of a function to linear and differential attacks is CCZ-invariant. CCZ-equivalent functions have also the same weakness/strength with respect to algebraic attacks. Indeed, let functions $`F,F^{}:๐ฝ_2^m๐ฝ_2^m`$ be CCZ-equivalent. Then $`F^{}=F_2F_1^1`$, where $`F_1,F_2`$ are defined by (1) and (2) for a certain linear permutation $`=(L_1,L_2)`$. If there exists a nonzero function $`\psi :๐ฝ_2^{2m}๐ฝ_2`$ of low degree such that
$$\psi (x,F(x))=0,x๐ฝ_2^m,$$
(3)
then $`\psi `$ could be used in an algebraic attack . The function $`\psi ^1`$ has the same degree as $`\psi `$ and Relation (3) is equivalent to
$$\psi ^1(F_1(x),F_2(x))=0,x๐ฝ_2^m,$$
which implies
$$\psi ^1(x,F^{}(x))=0,x๐ฝ_2^m$$
by replacing $`x`$ by $`F_1^1(x)`$. Hence, $`\psi ^1`$ could be used in an algebraic attack on $`F^{}`$, and vice versa. Therefore, the resistance of a function to algebraic attacks is also CCZ-invariant.
Since any permutation is CCZ-equivalent to its inverse then obviously the algebraic degree of a function is not CCZ-invariant (while it is EA-invariant as noted above). For example, if $`F:๐ฝ_2^m๐ฝ_2^m`$ is a Gold AB function then $`d^{}(F)=2`$ and $`d^{}(F^1)=\frac{m+1}{2}`$ as proven in .
EA-equivalent functions are CCZ-equivalent and if a function $`F`$ is a permutation then $`F`$ is CCZ-equivalent to $`F^1`$ . More precisely:
###### Proposition 3
Let $`F,F^{}:๐ฝ_2^m๐ฝ_2^m`$. The function $`F^{}`$ is EA-equivalent to the function $`F`$ or to the inverse of $`F`$ (if it exists) if and only if there exists a linear permutation $`=(L_1,L_2)`$ on $`๐ฝ_2^{2m}`$ such that $`(G_F)=G_F^{}`$ and the function $`L_1`$ depends only on one variable, i.e. $`L_1(x,y)=L(x)`$ or $`L_1(x,y)=L(y)`$.
*Proof.* Let $`(G_F)=G_F^{}`$ for some linear permutation $`=(L_1,L_2):๐ฝ_2^{2m}๐ฝ_2^{2m}`$ and $`L_1(x,y)=L(y)`$, $`L_2(x,y)=R_1(x)+R_2(y)`$, where $`L`$, $`R_1`$, $`R_2:๐ฝ_2^m๐ฝ_2^m`$ are linear. Then
$$F_1(x)=L_1(x,F(x))=LF(x),$$
$$F_2(x)=L_2(x,F(x))=R_1(x)+R_2F(x).$$
The function $`F_1`$ is a permutation, since $`(G_F)`$ is the graph of a function. Therefore, $`L`$ and $`F`$ have to be permutations. On the other hand, the system
$$\{\begin{array}{ccc}L(y)\hfill & =\hfill & 0\hfill \\ R_1(x)+R_2(y)\hfill & =\hfill & 0\hfill \end{array}$$
has only $`(0,0)`$ solution, since $`=(L_1,L_2)`$ is a permutation. But $`L`$ is also a permutation. Therefore, the linear function $`R_1`$ has to be a permutation too.
We have
$`F^{}(x)`$ $`=`$ $`F_2F_1^1(x)=R_1F^1L^1(x)+R_2FF^1L^1(x)`$
$`=`$ $`R_1F^1L^1(x)+R_2L^1(x).`$
Thus, $`F^{}`$ is EA-equivalent to $`F^1`$.
The proof that $`F^{}`$ is EA-equivalent to $`F`$, when $`L_1(x,y)=L(x)`$, is similar.
Conversely, let $`F^{}=R_1FR_2+R`$ or $`F^{}=R_1F^1R_2+R`$ for some linear permutations $`R_1,R_2`$ and for some linear function $`R`$. Then take $`(x,y)=(R_2^1(x),R_1(y)+RR_2^1(x))`$ in the first case and in the second case take $`(x,y)=(R_2^1(y),RR_2^1(y)+R_1(x))`$. Obviously, all conditions are satisfied. $`\mathrm{}`$
Remark 2 Proposition 3 shows that, for a function $`F:๐ฝ_2^m๐ฝ_2^m`$, if $`=(L_1,L_2)`$ and $`^{}=(L_1,L_2^{})`$ are linear permutations on $`๐ฝ_2^{2m}`$ such that the function $`L_1(x,F(x))`$ is a permutation, then the functions defined by the graphs $`(G_F)`$ and $`^{}(G_F)`$ are EA-equivalent. Indeed, we have $`(^{}^1)(G_F)=^{}(G_F)`$ and, denoting $`(^{}^1)=(S_1,S_2)`$, we have $`S_1(x,y)=x`$. This last equality can be easily checked by considering the linear functions $``$, $`^{}`$ and $`^1`$ as $`(2m)\times (2m)`$ matrices
$$=\left(\begin{array}{ccc}R_1\hfill & R_2\hfill & \\ T_1\hfill & T_2\hfill & \end{array}\right),^{}=\left(\begin{array}{ccc}R_1\hfill & R_2\hfill & \\ T_1^{}\hfill & T_2^{}\hfill & \end{array}\right),^1=\left(\begin{array}{ccc}A_1\hfill & A_2\hfill & \\ A_3\hfill & A_4\hfill & \end{array}\right).$$
We have
$$\left(\begin{array}{ccc}R_1\hfill & R_2\hfill & \\ T_1\hfill & T_2\hfill & \end{array}\right)\times \left(\begin{array}{ccc}A_1\hfill & A_2\hfill & \\ A_3\hfill & A_4\hfill & \end{array}\right)=\left(\begin{array}{ccc}R_1A_1+R_2A_3\hfill & R_1A_2+R_2A_4\hfill & \\ T_1A_1+T_2A_3\hfill & T_1A_2+T_2A_4\hfill & \end{array}\right)=\left(\begin{array}{ccc}I\hfill & 0\hfill & \\ 0\hfill & I\hfill & \end{array}\right),$$
where $`I`$ is the identity matrix and $`0`$ is the 0-matrix of order $`m`$, and this implies
$$\left(\begin{array}{ccc}R_1\hfill & R_2\hfill & \\ T_1^{}\hfill & T_2^{}\hfill & \end{array}\right)\times \left(\begin{array}{ccc}A_1\hfill & A_2\hfill & \\ A_3\hfill & A_4\hfill & \end{array}\right)=\left(\begin{array}{ccc}I& 0& \\ T_1^{}A_1+T_2^{}A_3& T_1^{}A_2+T_2^{}A_4& \end{array}\right).$$
$`\mathrm{}`$
Proposition 3 shows that if we want to construct functions $`F^{}`$ which are CCZ-equivalent to a function $`F`$ and EA-inequivalent to both $`F`$ and $`F^1`$ (if $`F^1`$ exists), then we have to use a linear function $`L_1(x,y)`$ depending on both variables. However, this condition is not sufficient as the following example shows.
Example 1 Let $`m=2n+1`$ and $`sn`$ \[mod $`2]`$. Then the linear function
$$(x,y)=(x+tr(x)+\underset{j=0}{\overset{ns}{}}y^{2^{2j+s}},y+tr(x))$$
is a permutation on $`๐ฝ_{2^m}^2`$ since the kernel of $``$ is $`\{(0,0)\}`$ (this can be checked by considering the cases $`tr(x)=0`$ and $`tr(x)=1`$). For the Gold AB function $`x^3`$ the function
$$F_1(x)=x+tr(x)+\underset{j=0}{\overset{ns}{}}(x^3)^{2^{2j+s}}$$
is a permutation on $`๐ฝ_{2^m}`$. Indeed, denoting $`L(y)=_{j=0}^{ns}y^{2^{2j+s}}`$ we have $`L(y+y^2)=y+tr(y)`$ and $`L((y+1)^3)=L(y^3)+y+tr(y)+1`$ since $`L(1)=1`$. Thus $`F_1(x)=L((x+1)^3)+1`$ and $`F_1`$ is a permutation if $`L`$ is bijective. The equation $`z=L(y)`$ implies $`z+z^2=y+tr(y)`$ and $`tr(z)=tr(y)`$. Therefore, $`L`$ is a permutation and $`L^1(x)=x+x^2+tr(x)`$, $`F_1^1(x)=[L^1(x+1)]^{\frac{1}{3}}+1`$. Finally, we get
$`F^{}(x)`$ $`=`$ $`F_2F_1^1(x)=([L^1(x+1)]^{\frac{1}{3}}+1)^3+tr([L^1(x+1)]^{\frac{1}{3}}+1)`$
$`=`$ $`L^1(x+1)+[L^1(x+1)]^{\frac{2}{3}}+[L^1(x+1)]^{\frac{1}{3}}+tr([L^1(x+1)]^{\frac{1}{3}})`$
$`=`$ $`L^1(x+1)+L^1([L^1(x+1)]^{\frac{1}{3}}).`$
Thus, both functions $`L_1`$ and $`L_2`$ depend on two variables, but the function $`F^{}`$ is EA-equivalent to the inverse of $`x^3`$.
This example can be generalized for any Gold AB function by replacing $`L^1(x)=x+x^{2^i}+tr(x)`$ and $`x^3`$ by $`x^{2^i+1}`$. $`\mathrm{}`$
If we want to classify all functions CCZ-equivalent to a given one $`F`$, then we should characterize all permutations of the form $`LF+L^{}`$, where $`L,L^{}`$ are linear. Indeed, we need to know which linear mapping $`L_1:๐ฝ_2^{2m}๐ฝ_2^m`$ is such that the function $`F_1(x)=L_1(x,F(x))`$ is a permutation. Clearly, $`L_1`$ can be written uniquely in the form $`L_1(x,y)=L(y)+L^{}(x)`$. If $`F_1`$ is a permutation then there exists a linear function $`L_2(x,y)`$ such that the linear function $`(L_1,L_2)(x,y)`$ is a permutation too. Indeed, $`L_1(x,F(x))`$ being a permutation, $`L_1`$ is onto and the kernel of $`L_1`$ has then dimension $`m`$. We can take for $`L_2`$ any linear permutation between $`Ker(L_1)`$ and $`๐ฝ_2^m`$, extended to $`๐ฝ_2^{2m}`$ by $`L_2(x+y)=L_2(x)`$ for all $`xKer(L_1)`$, $`yE`$, where $`E`$ is an $`m`$-dimensional subspace of $`F_2^{2m}`$ such that $`EKer(L_1)=๐ฝ_2^{2m}`$. Conversely, any linear function $`L_2`$ such that $`(L_1,L_2)`$ is a permutation has this form. Indeed, $`L_2`$ being onto, it has also an $`m`$-dimensional kernel, and $`(L_1,L_2)`$ being one to one, the kernels of $`L_1`$ and $`L_2`$ have trivial intersection. This proves that $`Ker(L_1)Ker(L_2)=๐ฝ_2^{2m}`$ and that we can take $`E=Ker(L_2)`$.
The following proposition gives a sufficient condition for a function to be EA-inequivalent to power functions.
###### Proposition 4
Let $`F`$ be a function from $`๐ฝ_{2^m}`$ to itself. If there exists an element $`c๐ฝ_{2^m}^{}`$ such that $`d^{}(tr(cF))\{0,1,d^{}(F)\}`$, then $`F`$ is EA-inequivalent to power functions.
*Proof.* It is proven in that for any power function $`x^d`$ and for any $`c๐ฝ_{2^m}`$, either the function $`tr(cx^d)`$ completely vanishes or it has exactly the algebraic degree $`w_2(d)`$. Thus, for any function $`F`$ which is affine equivalent to a power function, we have $`d^{}(tr(cF))\{0,d^{}(F)\}`$, $`c๐ฝ_{2^m}^{}`$. Therefore, if $`F`$ is EA-equivalent to a power function then $`d^{}(tr(cF))\{0,1,d^{}(F)\}`$, for every $`c๐ฝ_{2^m}^{}`$. $`\mathrm{}`$
## 4 CCZ-equivalence and the Gold functions
In the propositions below we give a characterization of permutations $`LF+L^{}`$ when $`F`$ is a Gold function. This characterization is not complete but it is useful for constructions of new APN and AB polynomials.
###### Proposition 5
Let $`F:๐ฝ_{2^m}๐ฝ_{2^m}`$, $`F(x)=L(x^{2^i+1})+L^{}(x)`$, where $`L,L^{}`$ are linear and $`\mathrm{gcd}(m,i)=1`$. Then $`F`$ is a permutation if and only if, for every $`u0`$ in $`๐ฝ_{2^m}`$ and every $`v`$ such that $`tr(v)=tr(1)`$, the condition $`L(u^{2^i+1}v)L^{}(u)`$ holds.
*Proof.* For any $`u๐ฝ_{2^m}^{}`$ we have
$`F(x)+F(x+u)=L(x^{2^i+1})+L^{}(x)+L((x+u)^{2^i+1})+L^{}(x+u)`$
$`=L\left(ux^{2^i}+u^{2^i}x+u^{2^i+1}\right)+L^{}(u)=L\left(u^{2^i+1}\left((x/u)^{2^i}+x/u+1\right)\right)+L^{}(u).`$
When $`x`$ ranges over $`๐ฝ_{2^m}`$ then $`(x/u)^{2^i}+x/u+1`$ ranges over $`\{v๐ฝ_{2^m}:tr(v)=tr(1)\}`$, since $`\mathrm{gcd}(m,i)=1`$. Hence, $`F`$ is injective (i.e. is a permutation) if and only if $`L(u^{2^i+1}v)L^{}(u)`$ for every $`u0`$ and every $`v`$ such that $`tr(v)=tr(1)`$. $`\mathrm{}`$
Remark 3 If $`m`$ is even, then, up to affine equivalence and without loss of generality, we can consider only functions of the type $`L(x^{2^i+1})+x`$. Indeed, if $`F(x)=L(x^{2^i+1})+L^{}(x)`$ is a permutation on $`๐ฝ_{2^m}`$ and $`m`$ is even, then it follows from Proposition 5 that $`L^{}`$ must be a permutation (take $`v=0`$). Then the function $`F^{}(x)=L^1F(x)=L^1L(x^{2^i+1})+x`$ is a permutation if and only if $`F`$ is a permutation. Moreover if a function $`L_2:๐ฝ_{2^m}^2๐ฝ_{2^m}`$ is such that the function $`(L(y)+L^{}(x),L_2(x,y))`$ is a permutation on $`๐ฝ_{2^m}^2`$, then obviously $`(L^1L(y)+x,L_2(x,y))`$ is a permutation too; note also that $`(L(y)+x,y)`$ is a permutation for any linear function $`L`$. Thus the function $`x^{2^i+1}`$ is CCZ-equivalent to the functions $`F_2F^1`$ and $`F_2F^1`$, where $`F_2(x)=L_2(x,x^{2^i+1})`$. We have $`F_2F^1(x)=F_2F^1L^{}(x).`$ Therefore, $`F_2F^1`$ and $`F_2F^1`$ are affine equivalent. $`\mathrm{}`$
###### Proposition 6
Let $`F:๐ฝ_{2^m}๐ฝ_{2^m}`$, $`F(x)=L(x^{2^i+1})+x`$, where $`L`$ is linear, $`m`$ even and $`\mathrm{gcd}(m,i)=1`$. Let $`L^{}`$ be the adjoint operator of $`L`$. Then $`F`$ is a permutation if and only if, for every $`v๐ฝ_{2^m}`$ such that $`L^{}(v)0`$, there exists $`u๐ฝ_{2^m}`$ such that $`L^{}(v)=u^{2^i+1}`$ and $`tr_{m/2}(\frac{v}{u})0`$, where $`tr_{m/2}`$ is the trace function from $`๐ฝ_{2^m}`$ to $`๐ฝ_{2^2}`$.
*Proof.* The function $`F`$ is a permutation if and only if, for every $`v0`$, the Boolean function $`tr(v(L(x^{2^i+1})+x))`$ is balanced (see e.g. ). Let $`L^{}`$ be the adjoint operator of $`L`$, then we have
$$tr(v(L(x^{2^i+1})+x))=tr(L^{}(v)x^{2^i+1}+vx).$$
If $`L^{}(v)=0`$, then the function $`tr(v(L(x^{2^i+1})+x))`$ is balanced. If $`L^{}(v)0`$, the quadratic function $`tr(L^{}(v)x^{2^i+1}+vx)`$ has associated symplectic form :
$$\phi (x,y)=tr(L^{}(v)x^{2^i}y+L^{}(v)xy^{2^i})=tr((L^{}(v)x^{2^i}+(L^{}(v)x)^{2^{mi}})y),$$
which has kernel :
$$=\{x๐ฝ_{2^m}:L^{}(v)x^{2^i}+(L^{}(v)x)^{2^{mi}}=0\}=$$
$$=\{x๐ฝ_{2^m}:L^{}(v)^{2^i}x^{2^{2i}}+L^{}(v)x=0\}=\{0\}\{x๐ฝ_{2^m}:L^{}(v)^{2^i1}x^{2^{2i}1}=1\}.$$
A quadratic function is balanced if and only if its restriction to the kernel of its associated symplectic form is not constant . The restriction of $`tr(L^{}(v)x^{2^i+1})`$ to $``$ is null. Indeed, $`L^{}(v)^{2^i1}x^{2^{2i}1}=1`$ implies that the order of $`L^{}(v)x^{2^i+1}`$ divides $`2^i1`$ and since $`\mathrm{gcd}(i,m)=1`$ then $`L^{}(v)x^{2^i+1}๐ฝ_2`$ and the trace function is null on $`๐ฝ_2`$, since $`m`$ is even. Hence, the function $`L(x^{2^i+1})+x`$ is a permutation if and only if every $`v`$ such that $`L^{}(v)0`$ lies outside the dual of $`\{0\}\{x๐ฝ_{2^m}:L^{}(v)^{2^i1}x^{2^{2i}1}=1\}`$. Equivalently, the function $`L(x^{2^i+1})+x`$ is a permutation if and only if, for every $`v`$ such that $`L^{}(v)0`$ the following two conditions hold:
1) $`L^{}(v)^{2^i1}`$ belongs to $`\{x^{2^{2i}1}:x๐ฝ_{2^m}\}`$ (otherwise, $``$ would be trivial), say $`L^{}(v)^{2^i1}=u^{2^{2i}1}`$, i.e. $`L^{}(v)=u^{2^i+1}`$ (since the mapping $`yy^{2^i1}`$ is a permutation); in this case $`=\{0\}\{x๐ฝ_{2^m}:(ux)^{2^{2i}1}=1\}=\frac{1}{u}\{y๐ฝ_{2^m}:y^{2^{2i}}=y\}=\frac{1}{u}๐ฝ_{2^j}`$, where $`j=\mathrm{gcd}(2i,m)=2`$, hence $`=\frac{1}{u}๐ฝ_4`$;
2) $`v`$ lies outside the dual of $``$, that is, $`L(vx)0`$ for some $`x`$ .
For every $`z๐ฝ_{2^m}`$ and every $`y๐ฝ_{2^2}`$ we have
$$tr(z\frac{1}{u}y)=tr_2(tr_{m/2}(\frac{z}{u}y))=tr_2(ytr_{m/2}(\frac{z}{u})).$$
Hence, the function $`L(x^{2^i+1})+x`$ is a permutation if and only if every $`v`$ such that $`L^{}(v)0`$ satisfies $`L^{}(v)=u^{2^i+1}`$ for some $`u`$ and $`tr_{m/2}(\frac{v}{u})0`$. $`\mathrm{}`$
## 5 New cases of AB and APN functions
The next theorems show that CCZ-equivalent functions are not necessarily EA-equivalent (including the equivalence to the inverse). They lead to infinite classes of almost bent and almost perfect nonlinear polynomials, which are EA-inequivalent to power functions.
For the function $`F^{}:๐ฝ_{2^m}๐ฝ_{2^m}`$, $`F(x)=x^{2^i+1}`$, $`\mathrm{gcd}(m,i)=1`$, and for any $`a๐ฝ_{2^m}^{}`$ the set
$$H_a=\{F(x+a)+F(x):x๐ฝ_{2^m}\}=\{\begin{array}{cc}\{x๐ฝ_{2^m}:tr(a^{(2^i+1)}x)=1\}\hfill & \text{if }m\text{ is odd}\hfill \\ \{x๐ฝ_{2^m}:tr(a^{(2^i+1)}x)=0\}\hfill & \text{if }m\text{ is even}\hfill \end{array}$$
is a linear hyperplane when $`m`$ is even and a complement of a linear hyperplane when $`m`$ is odd (see or Proposition 5).
We can use Proposition 1. For any $`a๐ฝ_{2^m}^{}`$ the set
$$V^{}=\{\begin{array}{cc}(0,๐ฝ_{2^m}\backslash H_a)(a,๐ฝ_{2^m}\backslash H_a)\hfill & \text{if }m\text{ is odd}\hfill \\ (0,H_a)(a,๐ฝ_{2^m}\backslash H_a)\hfill & \text{if }m\text{ is even}\hfill \end{array}$$
is a subgroup of $`๐ฝ_{2^m}^2`$ such that $`G_F`$ is transversal to $`V^{}`$ since $`A_FV^{}=\mathrm{}`$.
###### Theorem 1
The function $`F^{}:๐ฝ_{2^m}๐ฝ_{2^m}`$, $`m>3`$ odd,
$$F^{}(x)=x^{2^i+1}+(x^{2^i}+x)tr(x^{2^i+1}+x),\text{ }\mathrm{gcd}(m,i)=1,$$
is an AB function which is EA-inequivalent to any power function.
*Proof.* It is easy (but lengthy) to prove that for $`a๐ฝ_{2^m}^{}`$ the linear mapping
$$(x,y)=(x+atr(a^1x)+atr(a^{(2^i+1)}y),y+a^{2^i+1}tr(a^{(2^i+1)}y)+a^{2^i+1}tr(a^1x))$$
is an involution. We have
$$(0,y)=\{\begin{array}{cc}(a,y+a^{2^i+1})\hfill & \text{if }\text{ }yH_a\hfill \\ (0,y)\hfill & \text{if }\text{ }y๐ฝ_{2^m}\backslash H_a\hfill \end{array}.$$
Since $`tr(a^{(2^i+1)}a^{2^i+1})=1`$, then $`a^{2^i+1}H_a`$. The set $`H_a`$ is a complement of a linear hyperplane, hence the sum of any two elements from $`H_a`$ belongs to $`๐ฝ_{2^m}\backslash H_a`$. Therefore, $`^1(V)=(V)=V^{}`$ and by Proposition 1 the function $`F`$ is CCZ-equivalent to $`F_2F_1^1`$, where
$`F_1(x)`$ $`=`$ $`L_1(x,F(x))=x+atr(x/a)+atr((x/a)^{2^i+1}),`$
$`F_2(x)`$ $`=`$ $`L_2(x,F(x))=x^{2^i+1}+a^{2^i+1}tr((x/a)^{2^i+1})+a^{2^i+1}tr(x/a).`$
The function $`F_1`$ is an involution (this proves again that it is bijective):
$`F_1F_1(x)`$ $`=`$ $`x+atr(x/a)+atr((x/a)^{2^i+1})+atr(a^1(x+atr(x/a)`$
$`+`$ $`atr((x/a)^{2^i+1})))+atr(a^{(2^i+1)}(x+atr(x/a)+atr((x/a)^{2^i+1}))^{2^i+1})`$
$`=`$ $`x+3atr(x/a)+2atr((x/a)^{2^i+1})+atr(a^{(2^i+1)}(x^{2^i+1}+(ax^{2^i}+a^{2^i}x`$
$`+`$ $`a^{2^i+1})(tr(x/a)+tr((x/a)^{2^i+1}))))=x+atr(x/a)+atr((x/a)^{2^i+1})`$
$`+`$ $`atr((x/a)^{2^i}+(x/a)+1)(tr(x/a)+tr((x/a)^{2^i+1}))=x,`$
since
$$(x+atr(x/a)+atr((x/a)^{2^i+1}))^{2^i+1}=x^{2^i+1}+(ax^{2^i}+a^{2^i}x+a^{2^i+1})(tr(x/a)+tr((x/a)^{2^i+1}))$$
$$+2a^{2^i+1}tr(x/a)tr((x/a)^{2^i+1})=x^{2^i+1}+(ax^{2^i}+a^{2^i}x+a^{2^i+1})(tr(x/a)+tr((x/a)^{2^i+1})).$$
Therefore,
$`F_2F_1^1(x)`$ $`=`$ $`(x+atr(x/a)+atr((x/a)^{2^i+1}))^{2^i+1}+a^{2^i+1}tr(a^{(2^i+1)}(x+atr(x/a)`$
$`+`$ $`atr((x/a)^{2^i+1}))^{2^i+1})+a^{2^i+1}tr(a^1(x+atr(x/a)+atr((x/a)^{2^i+1})))`$
$`=`$ $`x^{2^i+1}+(ax^{2^i}+a^{2^i}x+a^{2^i+1})(tr(x/a)+tr((x/a)^{2^i+1}))+a^{2^i+1}tr((x/a)^{2^i}`$
$`+`$ $`(x/a)+1)(tr(x/a)+tr((x/a)^{2^i+1}))+2a^{2^i+1}tr(x/a)+2a^{2^i+1}tr((x/a)^{2^i+1})`$
$`=`$ $`a^{2^i+1}[(x/a)^{2^i+1}+((x/a)^{2^i}+(x/a))tr((x/a)+(x/a)^{2^i+1})]=a^{2^i+1}F^{}(x/a),`$
where $`F^{}(x)=x^{2^i+1}+(x^{2^i}+x)tr(x+x^{2^i+1})`$. By Proposition 2 the function $`F^{}`$ is AB.
It is not difficult to see that the algebraic degree of $`F^{}`$ is 3 for $`m>3`$. On the other hand $`tr(F^{}(x))=tr(x^{2^i+1})`$ and $`d^{}(tr(F^{}(x)))=2`$. It follows from Proposition 4 that $`F^{}`$ is EA-inequivalent to any power function. $`\mathrm{}`$
Remark 4 It was conjectured in that any AB function is EA-equivalent to a permutation. The AB function from Theorem 1 serves as a counterexample for this conjecture. It was checked by the help of a computer, that for no linear function $`L`$ on $`๐ฝ_{2^5}`$ the sum $`F^{}+L`$ is a permutation for the AB function $`F^{}(x)=x^{2^i+1}+(x^{2^i}+x)tr(x^{2^i+1}+x)`$, $`\mathrm{gcd}(5,i)=1`$. Thus, $`F^{}`$ is EA-inequivalent to any permutation but it is CCZ-equivalent to the permutation $`x^{2^i+1}`$. For $`m`$ even the existence of APN permutations on $`๐ฝ_{2^m}`$ is an open problem. It is shown by Nyberg that there exist no quadratic APN permutations when $`m`$ is even . But as we have seen the nonexistence of permutations EA-equivalent to the Gold functions does not mean yet that there exist no permutations CCZ-equivalent to them. $`\mathrm{}`$
###### Theorem 2
The function $`F^{}:๐ฝ_{2^m}๐ฝ_{2^m}`$, $`m4`$ even,
$$F^{}(x)=x^{2^i+1}+(x^{2^i}+x+1)tr(x^{2^i+1}),\text{ }\mathrm{gcd}(m,i)=1,$$
is an APN function which is EA-inequivalent to any power function.
*Proof.* For $`a๐ฝ_{2^m}^{},`$ the linear mapping
$$(x,y)=(L_1,L_2)(x,y)=(x+atr(a^{(2^i+1)}y),y)$$
is an involution and, obviously, $`(V)=V^{}`$. Thus, by Proposition 1 the function $`F`$ is CCZ-equivalent to $`F_2F_1^1`$, where
$`F_1(x)`$ $`=`$ $`L_1(x,F(x))=x+atr((x/a)^{2^i+1}),`$
$`F_2(x)`$ $`=`$ $`L_2(x,F(x))=x^{2^i+1}.`$
The function $`F_1`$ is an involution:
$`F_1F_1(x)`$ $`=`$ $`x+atr((x/a)^{2^i+1})+atr(a^{(2^i+1)}(x^{2^i+1}+ax^{2^i}tr((x/a)^{2^i+1})`$
$`+`$ $`a^{2^i}xtr((x/a)^{2^i+1})+a^{2^i+1}tr((x/a)^{2^i+1})))`$
$`=`$ $`x+2atr((x/a)^{2^i+1})+atr((x/a)^{2^i}+(x/a)+1)tr((x/a)^{2^i+1})=x.`$
We have
$`F_2F_1^1(x)`$ $`=`$ $`(x+atr((x/a)^{2^i+1}))^{2^i+1}=x^{2^i+1}+ax^{2^i}tr((x/a)^{2^i+1})`$
$`+`$ $`a^{2^i}xtr((x/a)^{2^i+1})+a^{2^i+1}tr((x/a)^{2^i+1})`$
$`=`$ $`a^{2^i+1}[(x/a)^{2^i+1}+((x/a)^{2^i}+(x/a)+1)tr((x/a)^{2^i+1})].`$
The function $`F_2F_1^1`$ is EA-equivalent to the function
$$F^{}(x)=x^{2^i+1}+(x^{2^i}+x+1)tr(x^{2^i+1}).$$
Hence, $`F^{}`$ is CCZ-equivalent to $`F`$ and it is APN by Proposition 2. The algebraic degree of $`F^{}`$ is $`3`$ and $`d^{}(tr(F^{}))=2`$ as $`tr(F^{}(x))=tr(x^{2^i+1})`$. Therefore, $`F^{}`$ is EA-inequivalent to power functions by Proposition 4. $`\mathrm{}`$
Note that the proofs of Theorems 1 and 2 do not depend on the condition $`\mathrm{gcd}(i,m)=1`$. When $`\mathrm{gcd}(i,m)=s`$ then the functions $`F^{}`$ have the same differential and linear properties as $`x^{2^i+1}`$ and, therefore, if $`m/s`$ is odd they can be considered as the first polynomials with three valued Walsh spectrum $`\{0,\pm 2^{\frac{m+s}{2}}\}`$, which are EA-inequivalent to power functions.
###### Theorem 3
The function $`F^{}:๐ฝ_{2^m}๐ฝ_{2^m}`$, $`m`$ divisible by $`6`$,
$$F^{}(x)=[x+tr_{m/3}(x^{2(2^i+1)}+x^{4(2^i+1)})+tr(x)tr_{m/3}(x^{2^i+1}+x^{2^{2i}(2^i+1)})]^{2^i+1},$$
with $`\mathrm{gcd}(m,i)=1`$, is an APN function.
*Proof.* The linear function $`:๐ฝ_{2^m}^2๐ฝ_{2^m}^2`$,
$$(x,y)=(L_1,L_2)(x,y)=(x+tr_{m/3}(y^2+y^4),y)$$
is a permutation since it has the kernel $`\{(0,0)\}`$. For $`m`$ divisible by $`6`$, the function
$$F_1(x)=L_1(x,F(x))=x+tr_{m/3}(x^{2(2^i+1)}+x^{4(2^i+1)})$$
is a permutation. Indeed, by Proposition 5 the function $`F_1`$ is a permutation if for every $`v๐ฝ_{2^m}`$ such that $`tr(v)=0`$ and every $`u๐ฝ_{2^m}^{}`$ the condition $`tr_{m/3}((u^{2^i+1}v)^2+(u^{2^i+1}v)^4)u`$ holds. Obviously, for any $`u๐ฝ_{2^3}`$ the condition is satisfied and when $`u๐ฝ_{2^3}^{}`$ the condition $`tr_{m/3}((u^{2^i+1}v)^2+(u^{2^i+1}v)^4)u`$ is equivalent to $`(u^{2^i+1}tr_{m/3}(v))^2+(u^{2^i+1}tr_{m/3}(v))^4u`$. Therefore, $`F_1`$ is a permutation if, for every $`u,w๐ฝ_{2^3}^{}`$, $`tr_3(w)=0`$ the condition $`(u^{2^i+1}w)^2+(u^{2^i+1}w)^4u`$ is satisfied and that was easily checked by a computer.
We show below that $`F_1^1=F_1F_1F_1F_1F_1`$.
We denote
$$T(x)=tr_{m/3}(x^{2^i+1}),$$
then
$$F_1(x)=x+T(x)^2+T(x)^4.$$
Since every element of $`๐ฝ_8`$ is equal to its $`8`$ power and the function $`tr_{m/3}(x)`$ is $`0`$ on $`๐ฝ_8`$, then
$`TF_1(x)=tr_{m/3}[(x+tr_{m/3}(x^{2(2^i+1)}+x^{4(2^i+1)}))^{2^i+1}]`$
$`=`$ $`tr_{m/3}[(x+tr_{m/3}(x^{2(2^i+1)}+x^{4(2^i+1)}))(x^{2^i}+tr_{m/3}(x^{2^{i+1}(2^i+1)}`$
$`+`$ $`x^{2^{i+2}(2^i+1)}))]=tr_{m/3}[x^{2^i+1}+xtr_{m/3}(x^{2^{i+1}(2^i+1)}+x^{2^{i+2}(2^i+1)})`$
$`+`$ $`x^{2^i}tr_{m/3}(x^{2(2^i+1)}+x^{4(2^i+1)})]=tr_{m/3}(x)tr_{m/3}(x^{2^{s+1}(2^i+1)}+x^{2^{s+2}(2^i+1)})`$
$`+`$ $`tr_{m/3}(x^{2^s})tr_{m/3}(x^{2(2^i+1)}+x^{4(2^i+1)})+tr_{m/3}(x^{2^i+1}),`$
where $`simod3`$.
Therefore,
$`F_1F_1(x)=F_1(x)+[TF_1(x))]^2+[TF_1(x))]^4`$
$`=`$ $`x+2tr_{m/3}(x^{2(2^i+1)}+x^{4(2^i+1)})+tr_{m/3}(x^2)tr_{m/3}(x^{2^{s+2}(2^i+1)}`$
$`+`$ $`x^{2^s(2^i+1)})+tr_{m/3}(x^{2^{s+1}})tr_{m/3}(x^{4(2^i+1)}+x^{2^i+1})`$
$`+`$ $`tr_{m/3}(x^4)tr_{m/3}(x^{2^s(2^i+1)}+x^{2^{s+1}(2^i+1)})+tr_{m/3}(x^{2^{s+2}})tr_{m/3}(x^{2^i+1}+x^{2(2^i+1)})`$
Considering separately the cases $`s=1`$ and $`s=2`$ we get
$`F_1F_1(x)`$ $`=`$ $`x+tr_{m/3}(x+x^2+x^4)tr_{m/3}(x^{2^i+1}+x^{2^s(2^i+1)})`$
$`=`$ $`x+tr(x)tr_{m/3}(x^{2^i+1}+x^{2^s(2^i+1)})=x+(T(x)+T(x)^{2^s})tr(x).`$
Like this we get
$$F_1F_1F_1F_1F_1(x)=x+T(x)^2+T(x)^4+tr(x)(T(x)+T(x)^{2^{2s}}),$$
$$F_1F_1F_1F_1F_1F_1(x)=x.$$
Thus,
$`F^{}(x)`$ $`=`$ $`F_2F_1^1(x)=[x+T(x)^2+T(x)^4+tr(x)(T(x)+T(x)^{2^{2s}})]^{2^i+1}=x^{2^i+1}+T(x)^{2^s+1}`$
$`+`$ $`tr(x^{2^i+1})T(x)^{2^{2s}}+tr(x)(T(x)+T(x)^4)+xtr(x)(T(x)+T(x)^{2^s})`$
$`+`$ $`x^{2^i}tr(x)(T(x)+T(x)^{2^{2s}})+x(T(x)+T(x)^{2^{2s}})+x^{2^i}(T(x)^2+T(x)^4),`$
where $`F_2(x)=L_2(x,F(x))=x^{2^i+1}.`$
The function $`F^{}`$ has the algebraic degree 4, this can be shown by lengthy but uncomplicated calculations. Hence, $`F^{}`$ is EA-inequivalent to other known APN functions since for $`m`$ divisible by 6 we have no known APN functions of algebraic degree 4. $`\mathrm{}`$
Let $`F:๐ฝ_{2^m}๐ฝ_{2^m}`$, $`F(x)=x^{2^i+1}`$ with $`\mathrm{gcd}(m,i)=1`$, $`m`$ odd and $`n`$ a divisor of $`m`$. Then $`A_FV^{}=\mathrm{}`$ for the subgroup
$$V^{}=\{(b,x):b๐ฝ_{2^n},x๐ฝ_{2^m},tr_{m/n}(x)=0\},$$
since if $`(b,x)V^{}`$ then $`tr(b^{(2^i+1)}x)=tr_n(b^{(2^i+1)}tr_{m/n}(x))=0`$ and $`(b,x)A_F`$. Hence, $`G_F`$ is transversal to $`V^{}`$. Using the subgroup $`V^{}`$ we construct an AB function $`F^{}`$ given in the following theorem.
###### Theorem 4
The function $`F^{}:๐ฝ_{2^m}๐ฝ_{2^m}`$, where $`m`$ is odd and divisible by $`n`$, $`mn`$ and $`\mathrm{gcd}(m,i)=1`$,
$`F^{}(x)`$ $`=`$ $`x^{2^i+1}+tr_{m/n}(x^{2^i+1})+x^{2^i}tr_{m/n}(x)+xtr_{m/n}(x)^{2^i}`$
$`+`$ $`[tr_{m/n}(x)^{2^i+1}+tr_{m/n}(x^{2^i+1})+tr_{m/n}(x)]^{\frac{1}{2^i+1}}(x^{2^i}+tr_{m/n}(x)^{2^i}+1)`$
$`+`$ $`[tr_{m/n}(x)^{2^i+1}+tr_{m/n}(x^{2^i+1})+tr_{m/n}(x)]^{\frac{2^i}{2^i+1}}(x+tr_{m/n}(x)),`$
is an AB function which is EA-inequivalent to any power function.
*Proof.* Let $`m`$ be odd and divisible by $`n`$. Obviously, the linear function
$$(x,y)=(L_1,L_2)(x,y)=(x+tr_{m/n}(x)+tr_{m/n}(y),y+tr_{m/n}(x))$$
is a permutation on $`๐ฝ_{2^m}^2`$ and
$$^1(x,y)=(x+tr_{m/n}(y),y+tr_{m/n}(x)+tr_{m/n}(y)).$$
We have
$$F_1(x)=x+tr_{m/n}(x)+tr_{m/n}(x^{2^i+1}),$$
$$F_2(x)=x^{2^i+1}+tr_{m/n}(x).$$
By Proposition 1, the function $`F_1`$ is a permutation since $`^1(V)=V^{}`$ and $`G_F`$ is transversal to $`V^{}`$. We need the inverse of the function $`F_1`$ to construct $`F^{}=F_2F_1^1`$.
For any fixed element $`x๐ฝ_{2^m}`$ we have
$$y=x+tr_{m/n}(x)+tr_{m/n}(x^{2^i+1})=x+u,$$
for some $`u๐ฝ_{2^n}`$, and, therefore, $`x=y+u`$. Then
$$y=(y+u)+tr_{m/n}(y+u)+tr_{m/n}((y+u)^{2^i+1})$$
which yields
$$u^{2^i+1}+u^{2^i}tr_{m/n}(y)+u(tr_{m/n}(y))^{2^i}+tr_{m/n}(y^{2^i+1})+tr_{m/n}(y)=0.$$
(4)
If $`tr_{m/n}(y)0`$ then we denote $`v=u/tr_{m/n}(y)`$ and we get
$$v^{2^i+1}+v^{2^i}+v+\frac{tr_{m/n}(y^{2^i+1})+tr_{m/n}(y)}{(tr_{m/n}(y))^{2^i+1}}=0.$$
Since $`v^{2^i+1}+v^{2^i}+v=(v+1)^{2^i+1}+1`$ then
$$v+1=\left[\frac{tr_{m/n}(y^{2^i+1})+tr_{m/n}(y)}{(tr_{m/n}(y))^{2^i+1}}+1\right]^{\frac{1}{2^i+1}}.$$
Replacing $`v`$ by $`u/tr_{m/n}(y)`$ we have
$$u=[(tr_{m/n}(y))^{2^i+1}+tr_{m/n}(y^{2^i+1})+tr_{m/n}(y)]^{\frac{1}{2^i+1}}+tr_{m/n}(y).$$
If $`tr_{m/n}(y)=0`$ then from the equality (4) we get $`u=[tr_{m/n}(y^{2^i+1})]^{\frac{1}{2^i+1}}`$ and we observe that $`u`$ equals again $`[(tr_{m/n}(y))^{2^i+1}+tr_{m/n}(y^{2^i+1})+tr_{m/n}(y)]^{\frac{1}{2^i+1}}+tr_{m/n}(y).`$ Thus, in all cases, we have
$$F_1^1(y)=y+u=y+[(tr_{m/n}(y))^{2^i+1}+tr_{m/n}(y^{2^i+1})+tr_{m/n}(y)]^{\frac{1}{2^i+1}}+tr_{m/n}(y)$$
and
$`F^{}(x)`$ $`=`$ $`F_2F_1^1(x)=[x+[(tr_{m/n}(x))^{2^i+1}+tr_{m/n}(x^{2^i+1})+tr_{m/n}(x)]^{\frac{1}{2^i+1}}+tr_{m/n}(x)]^{2^i+1}`$
$`+`$ $`tr_{m/n}[x+[(tr_{m/n}(x))^{2^i+1}+tr_{m/n}(x^{2^i+1})+tr_{m/n}(x)]^{\frac{1}{2^i+1}}+tr_{m/n}(x)]`$
$`=`$ $`x^{2^i+1}+tr_{m/n}(x^{2^i+1})+tr_{m/n}(x)+x^{2^i}tr_{m/n}(x)+x(tr_{m/n}(x))^{2^i}`$
$`+`$ $`[(tr_{m/n}(x))^{2^i+1}+tr_{m/n}(x^{2^i+1})+tr_{m/n}(x)]^{\frac{1}{2^i+1}}(x^{2^i}+(tr_{m/n}(x))^{2^i}+1)`$
$`+`$ $`[(tr_{m/n}(x))^{2^i+1}+tr_{m/n}(x^{2^i+1})+tr_{m/n}(x)]^{\frac{2^i}{2^i+1}}(x+tr_{m/n}(x)).`$
We show below that the function $`F^{}`$ has the algebraic degree $`n+2`$. It means that the number of functions CCZ-equivalent to a Gold AB function and EA-inequivalent to it is not smaller than the number of divisors of $`m`$.
The inverse of $`x^{2^i+1}`$ on $`F_{2^n}`$ is $`x^d`$, where
$$d=\underset{k=0}{\overset{\frac{n1}{2}}{}}2^{2ik},$$
and $`x^d`$ has the algebraic degree $`\frac{n+1}{2}`$ . Obviously, $`((tr_{m/n}(x))^{2^i+1}+tr_{m/n}(x^{2^i+1})+tr_{m/n}(x))^d`$ has the algebraic degree $`n+1`$ if and only if $`((tr_{m/n}(x))^{2^i+1}+tr_{m/n}(x^{2^i+1}))^d`$ has this algebraic degree.
We assume that $`mn`$ and $`n1`$ since when $`m=n`$ we get $`F^{}(x)=x^{\frac{1}{2^i+1}}+x`$ and Theorem 1 gives the case $`n=1`$.
We have
$$tr_{m/n}(x)=\underset{k=0}{\overset{\frac{m}{n}1}{}}x^{2^{kn}}$$
and
$$(tr_{m/n}(x))^{2^i+1}+tr_{m/n}(x^{2^i+1})=\underset{k=0}{\overset{\frac{m}{n}1}{}}x^{2^{kn}}\underset{k=0}{\overset{\frac{m}{n}1}{}}x^{2^{kn+i}}+\underset{k=0}{\overset{\frac{m}{n}1}{}}(x^{2^i+1})^{2^{kn}}$$
$$=\underset{k,j=0}{\overset{\frac{m}{n}1}{}}x^{2^{kn}+2^{jn+i}}+\underset{k=0}{\overset{\frac{m}{n}1}{}}x^{2^{kn}+2^{kn+i}}=\underset{\begin{array}{c}k,j=0\\ kj\end{array}}{\overset{\frac{m}{n}1}{}}x^{2^{kn}+2^{jn+i}}.$$
Note that we have
$$[(tr_{m/n}(x))^{2^i+1}+tr_{m/n}(x^{2^i+1})]^{2^{2i}+1}=\underset{\begin{array}{c}k,j=0\\ kj\end{array}}{\overset{\frac{m}{n}1}{}}x^{2^{kn}+2^{jn+i}}\underset{\begin{array}{c}k,j=0\\ kj\end{array}}{\overset{\frac{m}{n}1}{}}x^{2^{kn+2i}+2^{jn+3i}}$$
$$=\underset{\begin{array}{c}k,j,s,t=0\\ kj,st\end{array}}{\overset{\frac{m}{n}1}{}}x^{2^{kn}+2^{jn+i}+2^{sn+2i}+2^{tn+3i}}.$$
Similarly, we have
$$((tr_{m/n}(x))^{2^i+1}+tr_{m/n}(x^{2^i+1}))^d=\underset{(k_0,\mathrm{},k_n)I}{}x^{_{s=0}^n2^{k_sn+si}},$$
(5)
where $`I=\{(l_0,\mathrm{},l_n):0l_t\frac{m}{n}1,l_{2t}l_{2t+1}\}`$.
The equality $`k_sn+si=k_tn+ti`$ is possible for $`0s<tn`$ only for $`s=0`$ and $`t=n`$. Indeed, if $`k_sn+si=k_tn+ti`$ then $`(k_sk_t)n=(ts)i`$. Since $`\mathrm{gcd}(n,i)=1`$ and $`0t,sn`$ then $`t=n,`$ $`s=0`$ and $`k_sk_t=i`$.
For simplicity we consider now the equality (5) in the case $`i=1`$. In the sum $`_{s=0}^n2^{k_sn+s}`$ the largest possible item is $`2^{(\frac{m}{n}1)n+n}=2^m`$. Therefore, when $`k_n\frac{m}{n}1`$ the sum is smaller than $`2^m1`$. Besides, all items in the sum are different modulo $`2^m1`$ except the case when $`k_0=0`$ and $`k_n=\frac{m}{n}1`$ and in the cases where $`k_0=k_n+1`$. Therefore, when $`k_0=k_n=1`$ the number $`_{s=0}^n2^{k_sn+s}`$ has the weight $`n+1`$. On the other hand, when $`k_0=k_n=1`$ and $`k_1=k_{n1}=0`$ the term
$$x^{_{s=0}^n2^{k_sn+s}}$$
does not vanish in (5). Indeed, if
$$\underset{s=0}{\overset{n}{}}2^{k_sn+s}\underset{p=0}{\overset{n}{}}2^{t_pn+p}mod(2^m1)$$
then we have only two possibilities:
1) for any $`s`$ there exists $`p`$ such that $`k_sn+s=t_pn+p`$ (and vice versa). Then $`(k_st_p)n=ps`$ and since $`0s,pn`$ then $`k_0=t_n+1`$, $`t_0=k_n+1`$ and $`k_s=t_s`$ for $`s0,n`$. If $`k_0=k_n=1`$ then $`t_n=0`$, $`t_0=2`$. But in our case $`k_0=k_n=1`$, $`t_1=k_1=0`$, $`t_{n1}=k_{n1}=0`$ and, therefore, $`t_n0`$ since $`t_{n1}t_n`$.
2) if $`t_n=\frac{m}{n}1`$ then $`k_0`$ must be equal to $`0`$ or $`k_n=\frac{m}{n}1`$, but $`k_0=k_n=1`$.
Thus, when $`k_0=k_n=1`$ and $`k_1=k_{n1}=0`$ (for permissible $`k_s`$, $`1<s<n1`$) the term
$$x^{_{s=0}^n2^{k_sn+s}}$$
has the algebraic degree $`n+1`$ and it does not vanish in (5).
If $`n5`$ we can also take $`k_2=1,k_3=k_4=0`$ and then we get
$$\underset{s=0}{\overset{n}{}}2^{k_sn+s}=2^n+2+2^{n+2}+2^3+2^4+\mathrm{}+2^{2n}.$$
(6)
We have
$$((tr_{m/n}(x))^3+tr_{m/n}(x^3))^d(x^2+tr_{m/n}(x)^2)+((tr_{m/n}(x))^3+tr_{m/n}(x^3))^{2d}(x+tr_{m/n}(x))$$
$$=\underset{(k_0,\mathrm{},k_n)I}{}x^{_{s=0}^n2^{k_sn+s}}\underset{1km/n1}{}x^{2^{nk+1}}+\underset{(k_0,\mathrm{},k_n)I}{}x^{_{s=0}^n2^{k_sn+s+1}}\underset{1jm/n1}{}x^{2^{nj}}$$
$$=\underset{\begin{array}{c}(k_0,\mathrm{},k_n)I\\ 1km/n1\end{array}}{}x^{2^{nk+1}+_{s=0}^n2^{k_sn+s}}+\underset{\begin{array}{c}(k_0,\mathrm{},k_n)I\\ 1jm/n1\end{array}}{}x^{2^{nj}+_{s=0}^n2^{k_sn+s+1}}.$$
(7)
We consider the item with the exponent
$$2^n+2+2^{n+2}+2^3+2^4+\mathrm{}+2^{2n}+2^{nk+1}$$
(8)
from the first sum in (7). It is easy to see that $`2^n+2+2^{n+2}+2^3+2^4+\mathrm{}+2^{2n}+2^{nk+1}<2^m`$ since $`km/n1`$. In this sum $`nk+1=k_sn+s`$ only if $`s=1`$. But then $`k=k_1=0`$ which is in contradiction with $`1k`$. Thus, the number given by this sum has the weight $`n+2`$. The item with the exponent (8) does not vanish. Indeed, if there is another item in the first sum of (7) with this exponent then
$$2^n+2+2^{n+2}+2^3+2^4+\mathrm{}+2^{2n}+2^{nk+1}=2^{nj+1}+\underset{s=0}{\overset{n}{}}2^{k_sn+s}.$$
If $`k=j`$ then (6) is equal to another sum $`_{s=0}^n2^{k_sn+s}`$ and we already showed that it is impossible. If $`kj`$ then $`k_1=k`$ and $`j=0`$ while $`1j`$.
Assume there exists an item in the second sum of (7) with the exponent (8) then
$$2^n+2+2^{n+2}+2^3+2^4+\mathrm{}+2^{2n}+2^{nk+1}=2^{nj}+\underset{s=0}{\overset{n}{}}2^{k_sn+s+1}$$
for some $`j,`$ $`1jm/n1`$ and $`(k_0,\mathrm{},k_n)I`$. We have $`3=k_sn+s+1modm`$ for some $`s`$, $`0sn`$. Then $`k_sn=2s`$ or $`k_sn=m(s2)`$ and this is possible only if $`k_2=0`$ or $`k_2=m/n`$, but since $`0k_sm/n1`$, then $`3=k_sn+s+1modm`$ only if $`k_2=0`$. The same arguments show that $`4=k_sn+s+1modm`$ only if $`k_3=0`$ and that is in contradiction with the condition $`k_{2t}k_{2t+1}`$. Therefore the item with the exponent (8) does not vanish in (7) and then it does not vanish in the sum presenting the function $`F^{}`$. This completes the proof that $`F^{}`$ has the algebraic degree $`n+2`$.
The algebraic degree of the function $`tr(F^{}(x))`$ is not greater than $`n+1`$ since
$`tr(F^{}(x))`$ $`=`$ $`tr(x^{2^i+1}+tr_{m/n}(x^{2^i+1})+tr_{m/n}(x)+x^{2^i}tr_{m/n}(x)+x(tr_{m/n}(x))^{2^i})`$
$`+`$ $`tr_n([(tr_{m/n}(x))^{2^i+1}+tr_{m/n}(x^{2^i+1})+tr_{m/n}(x)]^{\frac{1}{2^i+1}}tr_{m/n}(x^{2^i}+(tr_{m/n}(x))^{2^i}+1))`$
$`+`$ $`tr_n([(tr_{m/n}(x))^{2^i+1}+tr_{m/n}(x^{2^i+1})+tr_{m/n}(x)]^{\frac{2^i}{2^i+1}}tr_{m/n}(x+tr_{m/n}(x)))`$
$`=`$ $`tr(x)+tr(x^{2^i}tr_{m/n}(x))+tr(x(tr_{m/n}(x))^{2^i})`$
$`+`$ $`tr([(tr_{m/n}(x))^{2^i+1}+tr_{m/n}(x^{2^i+1})+tr_{m/n}(x)]^{\frac{1}{2^i+1}}).`$
On the other hand $`d^{}(tr(F^{}(x)))`$ is not 0 or 1. Indeed, for any $`x๐ฝ_{2^n}`$ we have
$$tr(F^{}(x))=tr(x)+2tr(x^{2^i+1})+tr([2x^{2^i+1}+x]^{\frac{1}{2^i+1}})=tr(x)+tr(x^{\frac{1}{2^i+1}}).$$
The function $`tr(x^{\frac{1}{2^i+1}})`$ has the algebraic degree $`(n+1)/2`$. Indeed, $`d^{}(tr(x^{\frac{1}{2^i+1}}))\{0,d^{}(x^{\frac{1}{2^i+1}})\}`$ and since $`x^{\frac{1}{2^i+1}}`$ is a permutation then $`tr(x^{\frac{1}{2^i+1}})`$ is not a constant and $`d^{}(tr(x^{\frac{1}{2^i+1}}))=d^{}(x^{\frac{1}{2^i+1}})=(n+1)/2`$. Hence, $`tr(F^{}(x)`$ is not linear on $`๐ฝ_{2^n}`$ and then it cannot be linear on $`๐ฝ_{2^m}`$.
Thus, the function $`F^{}`$ is EA-inequivalent to any power function by Proposition 4. $`\mathrm{}`$
## 6 Conclusion
We have shown that there exist APN and AB functions which are EA-inequivalent to power functions, and therefore, which are new, up to EA-equivalence. But these functions are CCZ-equivalent to the Gold functions. We leave two open problems:
\- finding classes of functions CCZ-equivalent to other known APN or AB functions, which would be EA-inequivalent to all known APN and AB functions (or even, inequivalent to power functions);
\- finding classes of APN or AB functions which would be CCZ-inequivalent to all known APN and AB functions (or even, CCZ-inequivalent to power functions).
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# Nonequilibrium molecular dynamics of complex fluids near the gel point
## I Introduction
There have been many studies, both experimental and theoretical, of the rheological properties of complex fluids in the vicinity of the gel point adam96 . There is general agreement in the literature that the transition from the sol to the gel phase, at least in the case of chemical gels, is continuous and accompanied by the divergence of the shear viscosity $`\eta (p_cp)^s`$ where $`p`$ characterizes the degree of crosslinking or condensation and $`p_c`$ is the critical point. However, the experimental and theoretical values of the exponent $`s`$ are rather widely distributed and there is continuing debate concerning the existence of a single universality class for these transitions. We mention, in passing, that we are excluding the case of vulcanization from this discussion. There is substantial evidence, in the form of a Ginzburg criterion degen77 , that the vulcanization transition which involves the crosslinking of very long chains is for practical purposes mean-field-like.
There has been much less theoretical and experimental work on the normal stress differences. In a non-Newtonian fluid under shear flow, e.g., in a Couette geometry with the flow in the $`x`$-direction and velocity gradient in the $`z`$-direction, the first and second normal stress differences $`N_1=\sigma _{xx}\sigma _{zz}`$ and $`N_2=\sigma _{zz}\sigma _{yy}`$ are both not zero and, in the low shear-rate limit, proportional to the square of the shear rate $`\dot{\gamma }`$. It is conventional to define the normal stress coefficients $`\mathrm{\Psi }_1`$, $`\mathrm{\Psi }_2`$ through $`N_i=\mathrm{\Psi }_i\dot{\gamma }^2`$. One of the purposes of this work is to investigate the critical behavior of the normal stress coefficients. These are most easily calculated directly, i.e., by shearing the computational cell and measuring the normal stresses via a virial formula. While there exists a Green-Kubo formula broderix02 that yields the zero shear-rate limit of $`\mathrm{\Psi }_1`$, the integral involved presents significant computational problems close to the gel point: the decay of the relevant correlation function becomes extremely slow and truncation errors become unmanageable. For this reason, we have decided to utilize nonequilibrium molecular dynamics (NEMD) to calculate both the shear viscosity and the normal stress differences. We use the same model for which we have previously calculated the shear viscosity MP01 using the Green-Kubo formalism. Our NEMD results for the critical exponent $`s`$ are consistent with these earlier results. The NEMD simulations indicate that the first normal stress coefficient diverges at the critical point $`\mathrm{\Psi }_1(\dot{\gamma }=0)(p_cp)^\lambda `$ with $`\lambda 3.15`$, i.e., much more strongly than does the shear viscosity. The errors in $`\mathrm{\Psi }_2(\dot{\gamma }=0)`$ are much larger than those in $`\mathrm{\Psi }_1`$ and an independent determination of an exponent for its divergence is not feasible. However, the data are consistent with the conjecture that both normal stress coefficients diverge in the same way. We note that the conclusion that $`\mathrm{\Psi }_1`$ diverges more strongly than $`\eta `$ has also been arrived at by Broderix et al.broderix02 in the context of a Rouse-type model. However, their exponent $`\lambda 4.9`$ is significantly larger than ours.
The structure of this article is as follows. In section II we briefly describe the model that we have used and the computational techniques. Section III contains our results for the two-dimensional case and the results in three dimensions are presented in Section IV. We conclude in Section V with a brief discussion.
## II Model and Computational Methods
Our model of the sol phase is identical to that in MP01 . All particles interact through the soft sphere potential $`V(r_{ij})=ฯต(\sigma /r_{ij})^{36}`$ for $`r_{ij}1.5\sigma `$ and, in the three-dimensional calculations, we have used a single volume fraction $`\mathrm{\Phi }=\pi N\sigma ^3/6V=0.4`$ which is well below the liquid-solid coexistence density. All calculations were carried out at a temperature $`k_BT/ฯต=1`$. In the absence of crosslinks, this system is a simple liquid that has been well characterized powles . We initially placed the particles on the vertices of a simple cubic lattice and instantaneously and randomly introduced a fraction $`p`$ of nearest-neighbor bonds. We used the bonding potential $`V_b(r_{ij})=k(r_{ij}r_0)^2`$ where $`k=5ฯต/\sigma ^2`$ and where $`r_0=(\pi \mathrm{\Phi })^{1/3}/\sigma `$ so that there was no internal mechanical strain. This method of crosslinking ensures that the cluster size distribution is that of percolation in three dimensions and that a gel forms (in the thermodynamic limit $`N\mathrm{}`$) at $`p_c0.2488`$. Once the particles had been crosslinked they were free to move throughout the three-dimensional computational box. They were initially thermalized with periodic boundary conditions. Once equilibrium had been attained, the computational box was sheared at a rate $`\dot{\gamma }=v_x(z)/z`$ and the boundary conditions were changed to the Lees-Edwards boundary conditions allen . The system was then reequilibrated using the so-called SLLOD algorithm allen subject to the constraint that the kinetic energy in the frame following the overall flow of the particles remain constant. This kinetic energy is proportional to the square of the โpeculiar velocityโ $`(v_x\dot{\gamma }z,v_y,v_z)`$. Once a steady drift had been established, the diagonal elements of the stress tensor as well as $`\sigma _{xz}`$ were calculated from the appropriate virial formula. Calculations were performed for systems of $`N=L^3`$ particles with $`L`$ = 10, 15 and 20 over the entire range $`0p<p_c`$ and for shear rates generally in the range $`.005\dot{\gamma }0.1`$ in units of $`\sqrt{ฯต/m\sigma ^2}`$ although for a few cases larger shear rates were also imposed. It should be noted that if atomistic values of $`ฯต`$, $`m`$ and $`\sigma `$ are used, this range of shear rates corresponds to values of order $`10^{12}s^1`$, i.e., enormously large rates compared to experimental values. Even if one takes the point of view that the particles represent a colloidal suspension, the minimum shear rate is still of the order of $`10^3s^1`$. It was necessary to use such large shear rates in order to obtain reasonably well converged estimates for the normal stress coefficients, especially close to the gel point. For this range of shear rates, the equations of motion could be stably integrated with a time step $`\delta t=0.005\sqrt{m\sigma ^2/ฯต}`$ but for larger values of $`\dot{\gamma }`$ the time step had to be decreased.
The values of the shear viscosity, and to an even larger extent the normal stress coefficients, varied considerably for different realizations of the crosslinks. Therefore, we averaged the results over several thousand realizations even for values of $`p`$ as small as 0.1.
The same repulsive pair potential was used in the two-dimensional case. Here the particles were initially placed on a triangular lattice and instantaneously crosslinked as in three dimensions. This sets the gel point at $`p_c=2\mathrm{sin}(\pi /18)0.347296`$. The lattice constant of the initial configuration was $`1.2`$, and so these simulations were done at a number density $`n0.8`$. The spring constant used was $`k=40ฯต/\sigma ^2`$ and the value of $`r_0=1.2`$ was chosen to eliminate mechanical strain due to the crosslinks at zero temperature, as was done in three dimensions.
## III Results in Two Dimensions
The final results for the zero shear rate viscosity extrapolated from nonequilibrium molecular dynamics simulations are shown in Fig. 1. These data are from simulations of systems of size $`32\times 32=1024`$ particles. No significant deviation from a scaling form is visible for $`p>0`$ at this system size for the range of $`p`$ studied here. The power law divergence of the viscosity thus directly gives an estimate for the exponent $`s=2`$, without a finite size scaling analysis.
The inset to Fig. 1 shows the shear rate dependence of the viscosity as well as the zero shear rate values from the Green-Kubo formula, for an uncrosslinked and for a lightly crosslinked sample. The extrapolation of the finite shear rate values to zero shear-rates seems to be consistent with the Green-Kubo values. This fluid exhibits shear thinning, as is seen in some experiments on complex fluids. The determination of a zero-shear-rate viscosity requires fitting to some functional form for the shear-rate dependence of $`\eta `$. However, this value is not sensitive to the form chosen. Many different functional forms have been suggested for the shear-rate dependence of the viscosity, mostly suggested as phenomenological fitting functions to experimental data bhw . The value for $`\eta `$ in the main figure was estimated from a fit to a Lorentzian plus a constant term, as suggested in ferrario92 , while the difference between different fits was used as an estimate of the error. The Lorentzian form has the advantage that it is automatically symmetric in $`\dot{\gamma }`$ and is analytic near $`\dot{\gamma }=0`$.
There is only one normal stress difference, $`N_1=\sigma _{xx}\sigma _{yy}`$, in a two dimensional system. We have measured the associated coefficient $`\mathrm{\Psi }_1`$ close to $`p_c`$. This quantity is expected to diverge as a power law $`\mathrm{\Psi }_1(p_cp)^\lambda `$ as the gel transition is approached; the results of our simulation are shown in Fig. 2. We estimate an exponent $`\lambda 6`$ from our data.
The estimation of the normal stress coefficient was more difficult than for the viscosity. The division by $`\dot{\gamma }^2`$ produces large statistical errors for small $`\dot{\gamma }`$, thus making it difficult to determine the functional form of the shear-rate dependence. An example is shown in the inset to Fig. 2. Only a simple linear fit was used to estimate the zero-shear-rate normal stress coefficient.
## IV Results in Three Dimensions
In the three-dimensional case, we have results for systems of $`N=L^3`$ particles with $`L=10`$, $`15`$, and $`20`$. We first display, in Fig. 3, the shear-rate dependence of the viscosity $`\eta (\dot{\gamma })`$ for $`L=10`$ and five different values of the degree of crosslinking ranging from the simple fluid case, $`p=0`$, to a system close to the gel point ($`p=0.22`$). All systems show evidence of shear thinning, with this feature becoming much more prominent and setting in at lower shear rates as the gel point is approached. If we rescale the shear viscosity using the form $`\eta (p,\dot{\gamma })=a(p_cp)^s\stackrel{~}{\eta }(p,\stackrel{~}{\dot{\gamma }})`$ with $`s=0.7`$ and $`\stackrel{~}{\dot{\gamma }}=b(p_cp)^z\dot{\gamma }`$ where $`a`$ and $`b`$ are constants, we can achieve a respectable collapse of the data, as seen in Fig. 4. Since the data are noisy, we have not made a serious effort to optimize this collapse. Nevertheless, the โdynamicalโ exponent $`z2.35`$. The exponent $`s=0.7`$ used to rescale the viscosity is our best estimate of the exponent that governs the divergence of the zero shear-rate viscosity $`\eta (p,\dot{\gamma }=0)(p_cp)^s`$.
We note that dynamical scaling yields a connection between the exponent $`z`$ that controls the divergence of the longest relaxation time at the gel point and the exponents $`s`$ and $`t`$, i.e., $`z=s+t`$ adam96 . Here $`t`$ is the exponent that describes how the shear modulus vanishes as the gel point is approached from the solid side: $`\mu (pp_c)^t`$. For this model, we have determined mp99 that $`t2.0`$. This yields the prediction $`z2.7`$, a value not too far from the value used to rescale the shear rate.
The zero-shear-rate viscosity is shown in Fig. 5 in finite-size scaled form, i.e., we plot $`L^{s/\nu }\eta (p,L)`$ as function of $`(p_cp)^\nu L`$ where $`\nu `$ is the correlation length exponent of the three dimensional percolation problem $`\nu 0.88`$. As mentioned above, the value of the exponent $`s=0.7`$ is consistent with our previous result obtained from a Green-Kubo calculation MP01 .
We turn now to the normal stress coefficients. In Fig. 6 we plot the first normal stress coefficient $`\mathrm{\Psi }_1(p,\dot{\gamma })(\sigma _{xx}\sigma _{zz})/\dot{\gamma }^2`$ as function of the shear rate for $`L=10`$ and three values of $`p`$. For larger values of $`p`$, $`\mathrm{\Psi }_1`$ increases rapidly as the shear rate is decreased and an estimate of the zero shear-rate value $`\mathrm{\Psi }_1(p,\dot{\gamma })`$ is problematical. We have fit the data points to a second order polynomial $`\mathrm{\Psi }_1(p,\dot{\gamma })=\mathrm{\Psi }_1(p,0)+a\dot{\gamma }+b\dot{\gamma }^2`$ and obtained our estimate of the zero shear-rate value in this way. A fit to an exponential decay works equally well and produces estimates of $`\mathrm{\Psi }_1(p,0)`$ that differ by no more than 3% from those shown here. Similarly, the Lorentzian-plus-constant fit that was used to fit the viscosity in two dimensions also provides a reasonable fit to the data as long as there are enough values of the shear rate (more than five). The conclusions presented below are insensitive to the method of extrapolation.
The unscaled data for $`\mathrm{\Psi }_1(p,0)`$ are plotted in Fig. 7 for $`L=10`$, $`15`$ and $`20`$ as function of $`p_cp`$ along with a line representing the function $`a(p_cp)^\lambda `$ with $`\lambda =3.15`$ that captures the form of the data outside the critical region quite well. Closer to the critical point $`p_c0.2488`$, the usual finite-size effects that appear when the geometric correlation length $`\xi (p)`$ approaches the system size, $`L`$, are evident. These finite-size effects can be hidden by plotting $`L^{\lambda /\nu }\mathrm{\Psi }_1(p,L)`$ as a function of the scaled variable $`ฯตL/\xi L(p_cp)^\nu `$ where $`\nu 0.88`$ is the correlation length exponent. This is done in Fig. 8 and the data do collapse reasonably well to a universal curve. A useful consistency check on this procedure is available far from the critical point: if the exponents $`\lambda `$ and $`\nu `$ are correctly determined then the scaled normal stress coefficient should approach the power-law form $`ฯต^{\lambda /\nu }`$ at large $`ฯต`$. This line is also displayed in Fig. 8 and the data are consistent with the expected behavior.
We have also calculated the second normal stress coefficient $`\mathrm{\Psi }_2(p,L)`$. This coefficient is much smaller in magnitude than $`\mathrm{\Psi }_1`$ and negative, at least for our range of shear rates. The sample-to-sample fluctuations in $`N_2`$ are relatively much larger than those in $`N_1`$ and the data for $`\mathrm{\Psi }_2`$ are therefore much more noisy. Indeed, they are so poorly converged that a determination of a critical exponent from that dataset is not supportable. We have therefore carried out a rescaling of $`\mathrm{\Psi }_2`$ under the assumption that it diverges in the same way as $`\mathrm{\Psi }_1`$, i.e., controlled by the same critical exponent $`\lambda 3.15`$. The results of this rescaling are shown in Fig. 9 and we see that the data for $`L=15`$ and $`L=20`$ that are somewhat better converged than the data for $`L=10`$ seem to support such an assumption.
The authors of reference broderix02 have proposed a scaling relation for the exponent $`\lambda `$: $`\lambda =z+s`$ that seems to be rigorous for the Rouse model that they have used. In our case, using the two estimates of $`z`$ referred to above, namely $`z=2.35`$ and $`z=2.7`$, we obtain $`\lambda =3.05`$ and $`\lambda =3.4`$ which bracket the measured value. It should however be noted that in the Rouse model of broderix02 , the second normal stress coefficient $`\mathrm{\Psi }_2=0`$ for all values of $`p`$.
## V Conclusions
Using nonequilibrium molecular dynamics simulations, we have measured the divergence of both the viscosity and the normal stress coefficient in a model gel as the gel transition is approached. Our results for the divergence of the zero shear-rate viscosity are consistent with our previous calculation using a Green-Kubo formula to extract the viscosity from an equilibrium simulation. In addition, this model exhibits shear thinning as the shear rate is increased, as is observed in experiments on gelling systems. We found that the exponent governing the divergence of the shear viscosity to be $`s=0.7`$ in three dimensions. This value is consistent with some experiments devreux93 ; durand87 , as well as several analytical calculations muller03 . We also find $`s=2`$ in two dimensions.
We have also presented evidence that, in this model, the shear-rate-dependent viscosity can be rescaled onto a single universal curve. This indicates that the physics of the shear-thinning that we observe for all $`p`$ is the same close to the gel point as it is in the simple liquid ($`p=0`$).
The divergence of the normal stress close to the gel transition has not been measured in a experiment. As suggested in broderix02 , it should be possible to observe the very strong divergence of this quantity experimentally. Measuring both the viscosity and the normal stress close to the gel point would give the ratio of two dynamical exponents without determining the critical point, which is often difficult to determine accurately in an experiment. The ratio of $`s`$ to $`\lambda `$ would then provide a dynamical exponent which would characterize the universality class of a given material. Comparing this value to the values predicted by different models could then give some insight into which features of a microscopic model are important to the dynamical properties of an incipient gel.
## ACKNOWLEDGMENTS
We thank Bรฉla Joรณs for helpful discussions. This research was supported by the NSERC of Canada, and by the National Science Foundation through MRSEC grant No. DMR 0079909.
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# Superstellar clusters and their impact on their host galaxies.
## 1 The properties of superstellar clusters
The discovery by HST of a large population of unusually compact young superstellar clusters (SSCs) within starburst galaxies (see review by Ho 1997; Johnson et al. 2001; Colina et al. 2002; Larsen & Richtler 2000, Kobulnicky & Johnson 1999 and the proceedings edited by Lamers et al. 2004), has led us to infer the unit of massive or violent star formation. SSCs with masses in the range of a few $`\times 10^5`$ M to up to 6 $`\times 10^7`$ M (see Walcher et al. 2004; Pasquali et al. 2004) within a small volume of radius 3 to 10 pc, are indeed some of the most energetic entities found now in a large variety of galaxies. Note also that collections of them have now been found within a single starburst galaxy, as in M82 and the antennae galaxy (see Melo et al 2005; Whitmore et al. 1999).
The potential impact that these new units of star formation may have on the ISM of their host galaxies has been inferred from cluster synthesis models (see e.g. Cerviรฑo & Mas-Hesse 1994, Leitherer et al 1999). A coeval cluster with 10<sup>6</sup> M in stars, an initial mass function (IMF) similar to that proposed by Salpeter and an upper and lower mass range for the coeval event between 100 M and 1 M, leads initially to several tens of thousands of O stars. These however begin to disappear rather quickly (after t $``$ 3 Myr) as they complete their evolution and explode as supernovae (SNe). The cluster evolution is so rapid that after 10 Myr there are no O stars left within the cluster. All massive stars undergo strong stellar winds and all of them with a mass larger than 8 M will end their evolution exploding as supernova. Thus, one is to expect from our hypothetical cluster several tens of thousands of SNe over a time span of some 40 Myr. During the supernova phase a 10<sup>6</sup> M stellar cluster will produce an almost constant energy input rate of the order of 10<sup>40</sup> erg s<sup>-1</sup>. On the other hand, the ionizing luminosity emanating from the cluster would reach a constant value of 10<sup>53</sup> photons s<sup>-1</sup> during the first 3.5 Myr of evolution to then drastically drop (as t<sup>-5</sup>) as the most massive members of the association explode as supernova. The rapid drop in the ionizing photon flux implies that after 10 Myr of evolution, the $`UV`$ photon output would have fallen by more than two orders of magnitude from its initial value and the HII region that they may have originally produced would have drastically reduced its dimensions. Thus the HII region lifetime is restricted to the first 10 Myr of the evolution and is much shorter than the supernova phase. It is important to realize that only 10$`\%`$ of the stellar mass goes into stars with a mass larger than 10 M, however, it is this 10$`\%`$ the one that causes all the energetics from the starburst. Being massive, although smaller in number, massive stars also reinsert into the ISM, through their winds and SN explosions, almost 40$`\%`$ of the starburst total original mass. And thus from a starburst with an initial mass of 10<sup>6</sup> M one has to expect a total of almost 4 $`\times 10^5`$ M violently injected back into the ISM, during the 4 $`\times 10^7`$ years that the SN phase lasts. From these, almost 40,000 M will be in oxygen ions and less than 1000 M in iron (see Silich et al. 2001).
One of the features of the stellar synthesis models regarding the energetics of coeval star clusters is that they fortunately scale linearly with the mass of the starburst. It is therefore simple to derive the properties of starbursts of different masses, for as long as they present the IMF, metallicity and stars in the same mass range considered by the models.
When dealing with the outflows generated by star clusters, another important intrinsic property is the metallicity of their ejected matter. This is a strongly varying function of time, bound by the yields from massive stars and their evolution time. Thus, once the cluster IMF and the stellar mass limits are defined, the resultant metallicity of the ejected material is an invariant curve, independent of the cluster mass. Here we consider coeval clusters with a Salpeter IMF, and stars between 100 M and 1 M, as well as the evolutionary tracks with rotation of Maynet & Maeder (2002) and an instantaneous mixing of the recently processed metals with the stellar envelopes of the progenitors (see Silich et al. 2001 and Tenorio-Tagle et al. 2003, for an explicit description of the calculations). This leads to metallicity values (using oxygen as tracer) that rapidly reach 14 Z (see Figure 1), and although steadily decaying afterwards, the metallicity remains above solar values for a good deal of the evolution (for more than 20 Myr), to then fall to the original metallicity of the parental cloud, once the oxygen yield has been delivered. One of the main effects of an enhanced metallicity of the ejecta is to boost its radiative cooling and in such a case, massive clusters may inevitably enter into the strong cooling regime, to then find their stationary superwinds totally inhibited (see section 4).
## 2 Feedback from superstellar clusters
The close spacing between sources within a super-star cluster warrants a very efficient thermalization of all their winds and supernova explosions, leading to the high central overpressure that is to drive both a superbubble and in some cases a supergalactic wind (SGW). The outflow from the star cluster surface is fully defined by three quantities: The mass and mechanical energy deposition rates (hereafter $`\dot{M}_{SC}`$ and $`\dot{E}_{SC}`$) and the radius that encompasses the newly born sources ($`R_{SC}`$).
In the adiabatic case, the total mass and energy deposition rates define the temperature and thus the sound speed $`c_{SC}`$ at the cluster surface.
$$T_{SC}=\frac{0.3\mu }{k}\frac{\dot{E}_{SC}}{\dot{M}_{SC}},$$
(1)
where $`\mu `$ is the mean mass per particle and $`k`$ the Boltzmann constant. On the other hand, the density of matter streaming out of $`R_{SC}`$ is:
$$\rho =\frac{\dot{M}_{SC}}{4\pi R_{SC}^2c_{SC}},$$
(2)
Thus at $`R_{SC}`$ (see Chevalier & Clegg 1985; hereafter CC85), the ratio of thermal and kinetic energy flux to the total flux is
$$F_{th}/F_{tot}=\frac{\frac{1}{\gamma 1}\frac{P}{\rho }}{\frac{u^2}{2}+\frac{\gamma }{\gamma 1}\frac{P}{\rho }}=\frac{9}{20}$$
(3)
$$F_k/F_{tot}=\frac{u^2/2}{\frac{u^2}{2}+\frac{\gamma }{\gamma 1}\frac{P}{\rho }}=\frac{1}{4}.$$
(4)
There is however a rapid evolution as matter streams away from the central starburst. After crossing $`r=R_{SC}`$ the gas is immediately accelerated across the steep pressure gradients and rapidly reaches its terminal velocity ($`V_{\mathrm{}}2c_{SC}`$). This is due to a fast conversion of thermal energy, into kinetic energy of the resultant wind.
In a recent communication (Silich et al. 2003, 2004), we have revised the properties of SSCs by solving the flow equations without the assumption of an adiabatic flow made by Chevalier & Clegg (1985). In this case, the steady-state solution results from solving
$`{\displaystyle \frac{1}{r^2}}{\displaystyle \frac{\mathrm{d}}{\mathrm{d}r}}\left(\rho ur^2\right)=q_m,`$ (5)
$`\rho u{\displaystyle \frac{\mathrm{d}u}{\mathrm{d}r}}={\displaystyle \frac{\mathrm{d}P}{\mathrm{d}r}}q_mu`$ (6)
$`{\displaystyle \frac{1}{r^2}}{\displaystyle \frac{\mathrm{d}}{\mathrm{d}r}}\left[\rho ur^2\left({\displaystyle \frac{u^2}{2}}+{\displaystyle \frac{\gamma }{\gamma 1}}{\displaystyle \frac{P}{\rho }}\right)\right]=q_eQ,`$ (7)
where $`q_e`$ and $`q_m`$ are the energy and mass deposition rates per unit volume ($`q_e=q_m`$ = 0 if $`r`$ exceeds $`R_{SC}`$), $`Q`$ is the cooling rate ($`Q=n^2\mathrm{\Lambda }`$) where $`n`$ is the wind number density and $`\mathrm{\Lambda }`$ is the cooling function (a function of $`T`$ and metallicity). Central values of density and temperature are found by iteration, with the restriction that the flow velocity ought to increase from 0 km s<sup>-1</sup> at the cluster center to the sound speed when it reaches the cluster surface. If this happens then the stationary solution is fully warranted. A solution in which what is put in ($`\dot{M}_{SC}`$) is at all times equal to the outflow from the star cluster surface (4 $`\pi R_{SC}^2\rho _{SC}c_{SC}`$). Radiative cooling leaves almost unaffected the wind density and velocity distributions but its temperature may for energetic clusters fall rapidly to $`T10^4`$ K, close to the SSC surface. Figure 2 compares the adiabatic (dotted lines) with the radiative solution (solid lines) for a cluster with a $`\dot{E}_{SC}`$ = 3 $`\times 10^{41}`$ erg s<sup>-1</sup>, in which the wind temperature instead of steadily falling as $`r^{4/3}`$, it rapidly falls to $`T10^4`$ K, reducing drastically the size of the resultant X-ray emitting volume. The energetic winds lead to a total restructuring of the surrounding ISM, and in some extreme cases their energetics may even reach the outskirts of a galaxy causing a galactic wind.
We have also shown (see below section 4) that in the $`\dot{E}_{SC}`$ vs star cluster radius diagram, there is a threshold limit that massive and energetic compact clusters may cross to find out that radiative cooling inhibits their stationary outflow condition (Silich et al., 2004) and then the matter ejected by the stellar sources, unable to escape, is to accumulate within the star cluster volume (see section 4).
## 3 Negative feedback and the physics of supergalactic winds
In all cases below the threshold line, which may be considered quasi-adiabatic or strongly radiative, the energy dumped by the central starburst, is to cause a major impact on the surrounding gas. The supersonic stream leads immediately to a leading shock able to heat, accelerate and sweep all the overtaken material into a fast expanding shell. In this way, as the whole structure grows, the density, temperature and thermal pressure of the wind drops as $`r^2`$, $`r^{4/3}`$ and $`r^{10/3}`$, respectively (CC85).
Note however that such a flow is exposed to the appearance of reverse shocks whenever it meets an obstacle cloud or when its thermal pressure becomes lower than that of the surrounding gas, as it is the case in strongly radiative winds and within superbubbles. There, the high pressure acquired by the swept up ISM becomes larger than that of the freely expanding ejecta (the free wind region; FWR), where $`\rho `$, $`T`$ and $`P`$ are rapidly falling. The situation leads to the development of a reverse shock and with it to the thermalization of the wind kinetic energy, reducing the size of the free-wind region. Thus for the FWR to extend up to large distances away from the host galaxy, the shocks would have had to evolve and displace all the ISM, leading to a free path into the intergalactic medium through which the free wind may flow as a supergalactic wind. The energy required to achieve such a task, depends strongly on the ISM density distribution. As shown by Silich & Tenorio-Tagle (2001) the energy required to burst into the inter-galactic medium out of a fast rotating flattened galaxy is orders of magnitude smaller than that required to exceed the dimensions of a slow rotating and more spherical density distribution (see Figure 3).
Given their large UV photon output and mechanical energy input rate, SSCs are now believed to be the most powerful negative feedback agents in starburst galaxies, leading not only to a large-scale structuring of the ISM and to limit star formation, but also to be the agents capable of establishing as in M82 a supergalactic wind, thereby removing processed material from galaxies and causing the contamination of the IGM (Tenorio-Tagle et al. 2003).
### 3.1 The inner structure of M82
The biconical outflow of M82, the nearest example of a supergalactic wind (SGW), displays a collection of kpc long optical filaments embedded into an even more extended pool of soft X-ray emission. The outflow is known to extend even further, reaching the โH<sub>ฮฑ</sub> capโ at 11 kpc from the nucleus of M82 (see Devine & Bally 1999). Both features have been partly explained either with the results of Chevalier & Clegg (1985) model of an adiabatic, freely expanding, stationary wind and/or by the remnant of a large-scale superbubble evolving into the ISM and the halo of the galaxy (see for example Suchkov et al. 1994). Both explanations, based on the energetics of a single stellar cluster, fail to explain the detailed inner structure of the outflow. The elongated filaments for example, are now known to emanate from the central starburst and are not the result of a limb brightened superbubble outer structure (Ohyama et al 2002). Note also that the filaments cannot be reconcile with instabilities in the large-scale supershell, driven by matter entreinment, which in the models occurs at large distances, kpc from the energy source. Furthermore, the stationary superwind solution of Chevalier & Clegg leads to a laminar flow and not to gas condensation or to a filamentary structure at all. Also, as shown by Strickland & Stevens (2000) it has failed to matched the X-ray luminosity of the M82 superwind. Note also that the adiabatic assumption, central in the model of Chevalier & Clegg, has recently been shown to be inapplicable in the case of massive and concentrated starbursts (Silich et al. 2003, 2004). Another important issue not accounted by most of the numerical simulations is the size of the waist of the biconical structure (150 pc radius in the case of M82), which in all calculated cases under the assumption of a single source of energy, (perhaps with the only exception of Tenorio-Tagle & Muรฑoz-Tuรฑรณn 1997, 1998, which account for the infall of matter into the central starburst), also end up with a remnant that presents a wide open waist along the galaxy plane (see figures in Tomisaka & Ikeuchi, 1988; Suchkov et al., 1994). Further arguments regarding the disagreement between theory and observations are given in Strickland & Stevens 2000, Strickland, Ponman & Stevens 1997, and in Tenorio-Tagle et al. 2003.
### 3.2 The physics of supergalactic winds
So far, all calculations in the literature have assumed that the energy deposition arises from a single central cluster that spans several tens of pc, the typical size of a starburst. Following however, the indisputable observational findings with HST, we have recently made the first attempt to calculate the hydrodynamics that result from the interaction of the winds from neighboring young compact clusters present in a galaxy nucleus (see Tenorio-Tagle et al 2003). Several aspects were considered in our two dimensional approach to the problem. Among these, the metallicity of the superwind matter was shown to have a profound impact on the inner structure of supergalactic winds. Full three dimensional calculations of the interaction of multiple SSC winds are now underway.
Several two dimensional calculations using as initial condition CC85 adiabatic flows have been performed with the explicit Eulerian finite difference code described by Tenorio-Tagle & Muรฑoz-Tuรฑรณn (1997, 1998). This has been adapted to allow for the continuous injection of multiple winds (see below).
We have considered the winds from several identical SSCs, each with a mechanical energy deposition rate equal to 10<sup>40</sup> erg s<sup>-1</sup>. The energy is dumped at every time step within the central 5 pc of each of the sources following the adiabatic solution of Chevalier & Clegg (1985). The time dependent calculations do not consider thermal conductivity but do account for radiative cooling, with a cooling law (Raymond et al. 1976) scaled to the metallicity assumed for every case.
Figure 4 presents the results for which the assumed metallicity of the winds was set equal to 10Z, justified by the high metallicity outflows expected from massive bursts of star formation (see Figure 1). The winds from the SSCs are exposed to suffer multiple interactions with neighboring winds and are also exposed to radiative cooling. For the former, the issue is the separation between neighboring sources and for the latter the local values of density, temperature and metallicity. Radiative cooling would preferably impact the more powerful and more compact sources, leading to cold (T $``$ 10<sup>4</sup> K) highly supersonic streams.
Figure 4 shows three equally powerful ($`L_{SC}`$ = 10<sup>40</sup> erg s<sup>-1</sup>) superstellar clusters sitting at 0, 60 and 90 pc from the symmetry axis. All of them with an $`R_{SC}`$ = 5 pc, produce almost immediately a stream with a terminal velocity equal to 1000 km s<sup>-1</sup>. At t = 0 yr the three clusters are embedded in a uniform low density ($`\rho `$ = 10<sup>-26</sup> g cm<sup>-3</sup>) medium. Thus our calculations do not address the development of a superbubble, nor the phenomenon of breakout from a galaxy disk or the halo, into the IGM. The initial condition assumes that prior events have evacuated the region surrounding the superstellar clusters, and we have centered our attention on the interaction of the supersonic outflows.
Figure 4 shows the development of a SGW until it reaches dimensions of one kpc, together with the final temperature structure splitted into the four temperature regimes: The regime of H recombination 10<sup>4</sup> K - 10<sup>5</sup> K, followed by two regimes of soft X-ray emission 10<sup>5</sup> K - 10<sup>6</sup> K, and 10<sup>6</sup> K - 10<sup>7</sup> K and the hard X-ray emitting gas with temperatures between 10<sup>7</sup> K - 10<sup>8</sup> K.
The crowding of the isocontours in the figures indicates steep gradient both in density or in temperature and velocity, and thus traces the presence of shocks and of rapid cooling zones.
The interaction of neighboring supersonic winds causes the immediate formation of their respective reverse shocks, and of a high pressure region right behind them. The pressure (and temperature) reaches its largest values at the base of the interaction plane, exactly where the reverse shocks are perpendicular to the incoming streams. The high pressure gas then streams into lower pressure regions, defining together with radiative cooling, how broad or narrow the high pressure zones, behind the reverse shocks, are going to be.
This also happens if cooling is fast enough, the oblique reverse shocks rapidly acquire a standing location, however in these cases, the loss of temperature behind the shocks is compensated by gas condensation, leading to narrow, dense and cold filaments. The drastic drop in temperature occurs near the base of the outflow, where the gas density is large and radiative cooling is exacerbated. The dense structures are then launched at considerable speeds ($``$ several hundreds of km s<sup>-1</sup>) from zones near the plane of the galaxy. These dense and cold structures are easy target to the UV radiation produced by the superstellar clusters and thus upon cooling and recombination are likely to become photoionized. Note however that as the free winds continue to strike upon these structures, even at large distances from their origin, the resultant cold filaments give the appearance of being enveloped by soft X-ray emitting streams.
All of these shocks are largely oblique to the incoming streams and thus lead to two major effects: a) partial thermalization and b) collimation of the outflow. These effects result from the fact that only the component of the original isotropic outflow velocity perpendicular to the shocks is thermalized, while the parallel component is fully transmitted and thus causes the deflection of the outflow towards the shocks. This leads both, to an efficient collimation of the outflow in a general direction perpendicular to the plane of the galaxy, and to a substantial soft X-ray emission associated with the dense filamentary structure, extending up to large distances (kpc) from the plane of the galaxy. In the figures one can clearly appreciate that the oblique shocks, confronting the originally diverging flows, lead to distinct regions where the gas acquires very different temperatures, allowing for radiation in different energy bands.
From our results it is clear that a plethora of structure, both in X-rays and in the optical line regime, as in M82, may originate from the hydrodynamical interaction of neighboring winds. The interaction leads to multiple standing oblique (reverse) shocks and crossing shocks able to collimate the outflow away from the plane of the galaxy. In our two dimensional simulations, these are surfaces that become oblique to the diverging streams and thus evolve into oblique shocks that thermalize only partly the kinetic energy of the winds causing a substantial soft X-ray emission at large distances away from the galaxy plane. Surfaces that at the same time act as collimators, redirecting the winds in a direction perpendicular to the plane occupied by the collection of SSCs. Radiative cooling behind the oblique shocks leads, as soon as it sets in, to condensation of the shocked gas, and thus to the natural development of a network of filaments that forms near the base of the outflow, and streams away from the plane of the galaxy to reach kpc scales. Under many circumstances the filaments develop right at the base of the outflow and for all cases the prediction is that they are highly metallic. Hydrodynamic instabilities play also a major role on the filamentary structure. Nonlinear thin shell instabilities as studied by Vishniac (1994) as well as Kelvin Helmholtz instabilities, broaden, twist and generally shape the filaments as these stream upwards and reach kpc scales.
We thus postulate that if a collection of SSCs is sitting in a preferential plane, most of the injected energy would be channeled in a direction perpendicular to the plane of the host galaxy. This is achieved naturally as a consequence of the plethora of oblique and crossing shocks that redirect the initially isotropic winds. Collimation thus occurs without the need of a thick interstellar matter disk or a torus.
Our considerations point at a new set of possible parameters that profoundly impact the development of supergalactic winds. These are:
* The number and location of superstellar clusters within a galaxy nucleus.
* The intensity of star formation, or stellar mass in every superstellar cluster, which defines their mechanical luminosity,
* The age of the SSC, which impacts on the metallicity and thus on the cooling of the ejected matter.
All of these are relevant new parameters that may promote, as in M82, the inner structure of a supergalactic wind: the co-existense of X-rays and dense filaments, even at large distances from the sources of energy. Parameters that may promote self-collimation and with it the narrow waist of the biconical outflow. All of these features have been confirmed with full 3-D calculations, subject of a forthcoming communication. Note also that our results led to the full analysis of the HST data of M82 and to the discovery of 197 stellar clusters in the nucleus of M82 (see Melo et al. 2005).
## 4 Positive feedback from massive and compact SSC
The location of the threshold line in the $`\dot{E}_{SC}`$ vs size ($`R_{SC}`$) diagram (see Figure 5), the line that defines whether or not a wind is inhibited, depends on several variables. It depends on the size of the star-forming region ($`R_{SC}`$) and the metallicity of the ejected gas, which has a strong impact on the cooling curve. It also depends on the assumed $`\dot{E}_{SC}/\dot{M}_{SC}`$ or adiabatic terminal speed ($`v_{\mathrm{}}`$) of the wind. The latter is also bound to the usual assumption that the energy deposited by SN is always $`10^{51}`$ erg but the mass of the stars exploding within the cluster ranges from, say, 100 M to 8 M and so the injection speed (similar to $`v_{\mathrm{}}`$) and the deposited amount of matter are also functions of time. Another factor that strongly affects the location of the threshold line is the thermalization efficiency ($`ฯต`$) which defines the fraction of the mechanical energy that upon thermalization, can be evenly spread within the cluster volume. Estimates of $`ฯต`$ by several authors lead to values between 1 (Chevalier & Clegg; 1985) and 0.03 (see Melioli & Del Pino 2004 and references therein) and depends simply on the proximity of the sources undergoing winds and SNe, which through radiation may reduce the amount of energy available after thermalization. We have shown that there are three different types of solutions: SSCs far away from the threshold line (low mass, low energy clusters) undergo a quasi-adiabatic evolution well described by the Chevalier & Clegg (1985) solution. More energetic clusters are to have strongly radiative winds. Cooling hardly affects their velocity ($`v_wv_{\mathrm{}}`$) and density distribution ($`\rho _wr^2`$), but their temperature instead of falling as $`r^{4/3}`$, it falls rapidly to $`T_w10^4`$ K close to the SC boundary and the more so, the closer they are to the threshold line (see Figure 2). The strongly radiative winds around such clusters lead, compared to the adiabatic solution, to very much reduced X-ray envelope sizes. The third solution is for clusters above the threshold line. These would have their winds inhibited and as shown below, this turns them into very efficient positive feedback star-forming agents.
The facts above the threshold line are that radiative cooling drastically diminishes the sound speed $`c_{SC}`$ and the pressure gradient across the SSC volume, inhibiting the possibility of a wind. Radiative cooling upsets then the stationary condition in which the deposited matter ($`\dot{M}_{SC}`$) has to equal, at all times, the amount of matter streaming out of the SC volume
$$\dot{M}_{SC}=2\dot{E}_{SC}/v_{\mathrm{}}^2=4\pi R_{SC}^2\rho _{SC}c_{SC}$$
(8)
where $`v_{\mathrm{}}`$ is the resultant wind terminal speed in the absence of radiative cooling. As soon as this happens, the mass returned by the stars ($`\dot{M}_{SC}`$) begins to accumulate, promoting larger densities and an even faster cooling within the SSC volume. Following this trend, the accumulated gas density ultimately fulfill the gravitational instability criterion causing the collapse into a new stellar generation.
For clusters above the threshold line, Figure 6 shows how the density of the accumulating gas ($`\rho _{ac}=3\dot{M}_{SC}t/4\pi R_{SC}^3`$; where $`t`$ the evolution time), grows as a function of time within the SSC volume, until the moment when the density of the accumulating gas exceeds the gravitational instability criteria
$$\rho _J2.3\times 10^{20}\left(\frac{T}{100K}\right)\left(R_{SC1}\right)^2gcm^3$$
(9)
where $`\rho _J`$ is the Jeans density and $`R_{SC1}`$ is the SSC radius in pc units. At that moment, when $`\rho _{ac}=\rho _J`$ collapse will inevitably proceed. Note that for a given SSC (with a fixed volume), $`\rho _{Jeans}`$ is only a function of temperature and this is set, at least initially, by photoionization.
For massive clusters the initial ample supply of UV photons exceeds at first the number of recombinations within the volume occupied by the reinserted gas and the resultant HII region, given the large metallicities of the ejecta, is here assumed to rapidly approach an equilibrium temperature $`T_{HII}`$ a few $`10^3`$ K. At these temperatures the sound speed ($`<`$ 10 km s<sup>-1</sup>) remains well below the escape speed and the reinserted gas would inevitably continue to accumulate to rapidly (within $`1.5\times 10^6`$ yr) reach the value of the Jeans density for a gas at say, $`T_{HII}`$ = 5000 K, and collapse into a new stellar generation. The event gives rise to a new phase of matter accumulation, which once more will rapidly approach $`\rho _J`$ (for $`T`$= 5000 K) and undergo collapse within a free-fall time, of the order of 10<sup>5</sup> yr, while transforming $`2\times 10^5`$ M into stars. All stellar generations resultant from mass accumulation within the SSC volume have here been assumed to also acquire a Salpeter IMF with similar upper and lower mass limits as those imposed to the main superstellar cluster, and their resultant properties (mechanical energy and UV photon output) have been added to those produced by the main cluster.
A few, almost identical, stellar generations are expected from the accumulation process (solid rising lines in Figure 6), every time that the accumulated gas density $`\rho _{ac}`$ reaches $`\rho _J(5000K)`$ (dashed line in Figure 6). The situation changes slightly when the number of ionizing photons ($`N^0`$), despite the added contribution of secondary stellar generations, becomes insufficient to fully ionize the matter accumulated within the star cluster volume. This is due to the evolution of the main cluster, whose UV photon output begins to fall as $`t^5`$ after $``$ 3.5 Myr. Figure 6 shows $`\rho _{HII}`$ (thin solid line), the maximum density within the SSC volume that can be supported fully ionized by the UV radiation produced by the evolving cluster ($`\rho _{HII}=(3N^0\mu ^2/(4\pi R_{SC}^3\beta ))^{0.5}`$; where $`\mu =1.4m_H`$ and $`\beta `$, the recombination coefficient to to all levels above the ground level = 2.59 $`\times 10^{13}`$ cm<sup>-3</sup> s<sup>-1</sup>). During the accumulation process, once $`\rho _{ac}`$ exceeds $`\rho _{HII}`$, the ionized volume begins to shrink to end up as a collection of ultra compact HII regions around the most massive stars left within the cluster, while the bulk of the ejected material, now recombined, continues to cool, approaching rapidly a temperature $``$ 100 K. The characteristic cooling time ($`t_\mathrm{\Lambda }=3kT/2\mathrm{\Lambda }n`$; where $`\mathrm{\Lambda }`$ is the cooling rate, a function of T and Z) is also very small, $``$ 1.5 $`\times 10^5`$ yr, and is to become even shorter as matter continues to accumulate. Matter is at all times uniformly replenished within the whole SSC volume, and thus the gas density presents an almost uniform value. However, the accumulating gas now has two different temperatures ($`T_{HII}`$ and 100 K) and as $`\rho _{ac}`$ grows and the fraction of the ionized volume ($`f_{HII}=(3N^0\mu ^2)/(4\pi R_{SC}^3\beta \rho _{ac}^2`$)) shrinks, the size of cold condensations (at 100 K) able to become gravitationally unstable and their free-fall time also become smaller.
The drop in the number of ionizing photons and the consequent growth of the neutral volume, lead then to a second important condition in which the characteristic accumulation time
$$\tau _{ac}=\frac{4}{3}\pi \rho _{gas}R_{SC}^3(1f_{HII})/\dot{M}$$
(10)
becomes equal to the free-fall time
$$t_{ff}=\sqrt{\frac{3\pi }{32G\rho _{gas}}},$$
(11)
This condition defines $`\rho _{gas}`$, the density above the Jeans instability limit for a neutral condensation at 100 K ($`\rho _J`$(100K)):
$$\rho _{gas}=\left(\frac{27\dot{M}^2}{512\pi G\left(1f_{HII}\right)}\right)^{1/3}R_{SC}^2=1.15\times 10^{18}\left(\frac{\dot{M}}{1M_{}yr^1}\right)^{2/3}\left(\frac{R_{SC}}{1pc}\right)^2$$
(12)
and thus once $`\rho _{ac}`$ becomes equal to $`\rho _{gas}`$, a new stationary solution becomes possible.
Everything happens very rapidly, compared to the evolution time-scale of the parental cluster ($``$ 40 Myr), and almost at the same time. The ejected matter is thermalized within the SSC volume and immediately begins to cool. At the same time that it accumulates making cooling even faster. This allows it to rapidly reach the required $`\rho _{gas}`$ value, above the Jeans instability limit, that warrants its collapse in a similar time-scale, while the collapsing material is replenished by the newly ejected matter. When this happens, the mass deposition rate from the cluster becomes equal to the rate of star formation. Gravitational collapse and star formation within the star cluster volume and with the matter injected by all sources, drive in this way a new stationary condition through a new era of quasi-continuous star formation in which $`\dot{M}_{SC}`$ is now equal to the star formation rate (SFR). Further details describing this phase are given in Tenorio-Tagle et al. 2005.
## 5 Conclusions
Superstellar clusters are certainly the main mode of massive star formation in starburst and interacting galaxies. We have reviewed here how is that they work and the possible impact that they may have into the surrounding ISM. We have revised the assumptions of Chevalier & Clegg (1985) dealing with thermalization and the flow requirements to establish a stationary superwind emanating from these sources. We have also emphasized the importance of radiative cooling and how does it affect the superwind X-ray envelopes.
Calculations in the literature have left clear the fact that single energy sources lead to superbubbles and supershells. However, to produce a supergalactic wind with a detailed inner structure as in M82, multiple sources, seem to be required.
We have also shown that a straight forward extrapolation towards the most massive and compact coeval clusters is not valid. When radiative cooling becomes significant within the SSC volume itself, then instead of driving a superwind able to disperse the surrounding ISM and even channel its way into the IGM, events that have make them been regarded as negative feedback agents, they become in fact extreme examples of positive star formation feedback.
Massive and compact coeval clusters appear in the $`\dot{E}_{SC}`$ vs cluster size diagram above the threshold line, in the region where radiative cooling inhibits the development of stationary superwinds. In such cases the matter reinserted, through stellar winds and supernovae, is unable to escape and after a short phase of matter accumulation, a new stationary solution in which $`\dot{M}_{SC}`$ becomes equal to the SFR is rapidly met. A positive feedback condition in which new stellar generations result in situ, from the collapse of the matter reinserted by the sources evolving within the star cluster volume.
The massive concentrations imply a high efficiency of star formation which permits even after long evolutionary times the tight configuration that characterizes them, despite stellar evolution and its impact through photo-ionization, winds and supernova explosions, believed to efficiently disperse the gas left over from star formation. It is thus the self-gravity that results from the high efficiency what keeps the sources bound together.
The secondary star formation process while causing a faster mass deposition rate, drives the SFR to grow from 0.1 to 0.25 M yr<sup>-1</sup> over the parent cluster supernova phase ($``$ 40 Myr). The continuous reprocessing of the ejected material leads effectively to a continuous transformation of the high mass stars into a low mass ($``$ 8 M) population, keeping constant the total mass of the stellar component.
A central issue, regarding ISM studies, is the fact that the more massive and compact clusters (as those detected by HST and other large ground-based telescopes), are unable to generate superwinds and shed their metals into the ISM or the IGM. Their evolution leads to many stellar generations and thus to a mixture of stellar populations, all contaminated by the products from former stellar generations. An exacerbated episode of star formation that leaves no trace of its evolution in the ISM.
GTT acknowledges finantial support from the Secretarรญa de Estado de Universidades e Investigaciรณn (Espaรฑa) ref: SAB2004-0189 and the hospitality of the Instituto de Astrofรญsica de Andalucรญa (IAA, CSIC) in Granada, Spain. This study has been partly supported by AYA2004-08260-CO3-O1 from the Spanish Consejo Superior de Investigaciones Cientรญficas.
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# Rรฉfรฉrences
Geometry of certain Lie-Frobenius groups
Jean Michel Dardiรฉ, Alberto Medina and Hassรจne Siby
Abstract. Let be $`G_{n,p}(๐)=M_{n,p}(๐)GL(๐^n)`$, $`๐=`$ or $``$, the semi-direct product of the additive group of matrices $`M_{n,p}(๐)`$ by the group $`GL(๐^n)`$ where the action is done by multiplication of matrices. If $`n=kp`$,$`p1`$ the Lie group $`G_{n,p}(๐)`$ admits an exact left invariant symplectic form . We study the geometry of this symplectic manifold. If $`n>p`$ ( resp. $`n=p`$ ) we prove that $`G_{np,p}(๐)`$ (resp $`G_{n\mathrm{1,1}}(๐)`$ ) is the symplectic reduction of the symplectic orthogonal $`(M_{n,p}(๐))^{}`$ of $`M_{n,p}(๐)`$ in $`G_{n,p}(๐)`$ and reciprocally that $`G_{n,p}(๐)`$ is a symplectic double extension, in the sense of , of $`G_{np,p}(๐)`$ ( resp. of $`G_{n\mathrm{1,1}}(๐)`$). Moreover we show that $`G_{n,p}(๐)`$ admits two left transverse Lagrangian foliations (with affine and closed leaves) . Consequently there exists (a left invariant) a canonical torsion free symplectic connection on $`G_{n,p}(๐)`$.
Key words:
Exact symplectic Lie group , Affine Lie group, Symplectic Reduction , Symplectic double extension.
Introduction The group $`G_{n\mathrm{,1}}(๐)=Aff(๐^n)`$ where $`๐=`$ or $``$ admits left invariant symplectic structures which are all exact because $`H^2(aff(๐^n),๐)=0`$().
To give an invariant symplectic form on $`G_{n\mathrm{,1}}(๐)`$ is to give a linear form $`\alpha `$ on the Lie algebra $`aff(๐^n)`$ such that the 2-cobound $`\delta \alpha `$ is non degenerate. That is to say $`\alpha `$ is a point with an open orbit under the coadjoint action of $`G_{n\mathrm{,1}}(๐)`$. These points have been characterized in within the framework of study of affine group representations and in where is carried on the study of left invariant symplectic structures of affine group initiated in . In , it showed that coadjoint action of Lie group semi-direct product of $`M_{n,p}(๐)`$ and $`GL(๐^n)`$ by matricial product, where $`M_{n,p}(๐)`$ indicate the additive group of $`(n\times p)`$-matrices, admits open orbits if only if $`n=kp`$. We extend to these groups denoted $`G_{n,p}(๐)`$ the result obtained in for the classical affine group: existence of a unique symplectic structure up isomorphism (Theorem 2.7), existence of transverse symplectic foliations (Theorem 2.6), as well as of a Lagrangian bi-foliation with closed leaves (Theorem 3.2). Such a pair of Lagrangian foliations is important in the quantization procedure (polarization) and implies the existence of a canonical (torsion free) symplectic connection on $`G_{n,p}(๐)`$. The natural left action of $`G_{n,p}(๐)`$ on itself being hamiltonian we can provide (Theorem 2.1) a similar of structureโs theorem of simply connected symplectic Lie group from .
In the third part we study up the fibrations of this theorem for explicit suggest a construction of $`(G_{n,p}(๐),d\alpha _1^+)`$ when the generalized double extension of donโt apply in this case.
To complete this introduction recall that a symplectic Lie group $`(G,\omega ^+)`$ has an affine structure given by the left invariant connection ( torsion free and zero curvature ) $``$ defined for all $`x,y,zT_\epsilon (G)`$ by :
$$\omega ^+(_{x^+}^+y^+,z^+)=\omega ^+(y^+,[x^+,z^+])$$
where $`x^+`$ denotes the left invariant vector field associated to $`x`$.
Then one says that the pair $`(G,^+)`$ is an affine Lie group. In this case the product on $`๐ข`$, $`xy=(_{x^+}^+y^+)`$ is with associator left invariant and verifies the condition $`xyyx=[x,y]`$. This structure plays a central rule in this study. 1. Left invariant symplectic structures on the Lie groups $`G_{n,p}`$ In what follows $`G_n(๐)`$ denotes the classical affine group where $`๐=`$ or $``$ and $`๐ข_n(๐)`$ is its Lie algebra. By analogy, we denote $`G_{n,p}(๐)`$ the group semi-direct product $`M_{n,p}(๐)Gl(๐^n)`$ with $`n=kp,p1`$ and $`๐ข_{n,p}(๐^n)`$, its Lie algebra. Obviously $`G_{n\mathrm{,1}}(๐)`$ is isomorphic to $`Aff(๐^n)`$. Considering that $`H^2(aff(๐^n),๐)=0`$ (see ) any left invariant alternate 2-form is invariantly exact. This result becomes general to $`๐ข_{n,p}(๐)`$ in the following way:
1.1 Lemma. Every left invariant symplectic form on $`G_{n,p}(๐)`$ is exact.
Proof. If $`k=p=1`$; $`๐ข_{\mathrm{1,1}}`$ is isomorphic to $`aff(๐)`$ and the result follows.
Assume that $`k`$ (or $`p`$) is greater than 1.Let be $`\omega Z^2(๐ข_{n,p};๐)`$, $`๐ข_{n,p}=M_{p,n}(๐)gl(n)`$. This means that we have
$$\omega ([a,b],c)=0(\text{*})\text{for every }a\text{,}b\text{,}c\text{ in }๐ข_{n,p}.$$
If we take $`a=(x,u)`$, $`b=(y,v)`$ and $`c=(0,Id)`$
(*) implies $`\omega (x,y)=0`$ for every $`x,yM_{n,p}`$, and consequently $`\omega (x,b)+\omega (a,y)=\omega (I,uyvx+[u,v])`$.
But this is equivalent to
$$\omega (a,b)=\omega (I,uxvy)$$
$$=\delta \beta (a,b)\text{ with }\beta =\omega (I,)\text{ }\mathrm{}$$
N.B In fact for all Lie algebra $`๐ข`$ having an element $`a๐ข`$ such that $`ad_a`$ is projector, we have $`H^2(๐ข,๐)=0`$.
Remark 1. The mapping
$$M_{p,n}(๐)\times gl(n)(M_{n,p}(๐)\times gl(n))^{}$$
$$(H,M)\alpha _{(H,M)}$$
given by
$$\alpha _{(H,M)}(N,V)=tr(N.H)+tr(M.V)$$
is a linear isomorphism.
In the following, the dual space $`๐ข_{n,p}^{}`$ of $`๐ข_{n,p}`$ is identified with $`M_{p,n}(๐)\times gl(n)`$. To give an invariant closed $`2`$-form on $`G_{n,p}(๐)`$ it is equivalent to give a linear form $`\alpha `$ on $`๐ข_{n,p}(๐)`$ . The $`2`$-form $`d\alpha ^+`$, where $`\alpha ^+`$ denotes the left invariant $`1`$-form defined by $`\alpha `$ is symplectic if and only if $`\alpha `$ has an open orbit under the coadjoint representation. Using the Remark 1 and the natural embedding of $`๐ข_{n,p}`$ in $`GL(๐^{n+p})`$ we can show 1.2 Lemma. The coadjoint representations of $`๐ข_{n,p}`$ and $`G_{n,p}`$ are given by the following formulas:
$`(i)`$
$$ad_{(0,u)}^{}(H,N)=(H,[u,N])$$
$`(ii)`$
$$ad_{(x\mathrm{,0})}^{}(H,N)=(0,x.H)$$
$`(iii)`$
$$Ad_{(0,U)}^{}(H,N)=(H.U^1,UNU^1)$$
$`(iv)`$
$$Ad_{(X,Id_{๐^n})}^{}(H,N)=(H,N+X.H)$$
This implies
$$Ad_{(X,U)}^{}(H,N)=(H.U^1,UNU^1+XHU^1)$$
Consequently the coadjoint orbit $`Ad_{G_{n,p}}^{}(\alpha )`$ of $`\alpha =(H_0,N_0)`$ with $`H_0=(0,I_p)M_{p,n}`$ and
$$N_0=\left(\begin{array}{ccc}0& 0& 0\\ I_p& 0& 0\\ 0& I_p& 0\end{array}\right)$$
is open. Denote by $`H_0`$ the $`(p,n`$)-matrix whose $`p\times p`$ blocks are all null except for the last, which is the identity of $`๐^p`$ and denoted by $`N_0`$ the $`(n\times n)`$-matrix whose $`p\times p`$ blocks are all null except the sub-diagonals which are $`Id_{๐^p}`$. The previous formulas allow us to verify that the orbit of $`(H_0,N_0)`$ is open since the isotropy subalgebra is trivial. 2. Symplectic Reduction - Left invariant Symplectic foliation on $`G_{n,p}(๐)`$.
Denote by $`\omega ^+=d\alpha ^+`$, $`\alpha =(H,N)๐ข_{n,p}^{}`$ a left invariant symplectic form on $`G_{n,p}`$.
The action of $`G_{n,p}(๐)`$ on $`G_{n,p}(๐)`$ given by
$$L_G:G_{n,p}(๐)\times G_{n,p}(๐)G_{n,p}(๐),(\tau ,\sigma )\tau \sigma (\text{ product in }G_{n,p}\text{ })$$
is a symplectic action. Morever $`L_G`$ is a Hamiltonian action; a momentum mapping for the action is given by
$$\mu :G_{n,p}(๐)๐ข_{n,p}(๐)^{}$$
$$\sigma \mu (\sigma ):x\alpha ^+(\sigma ),x^{}(\sigma )$$
where $`x๐ข_{n,p}`$ and $`x^{}`$ is the right invariant vector field associated to $`x`$.
The subgroup $`:=M_{n,p}`$ is (totally) isotropic for $`\omega ^+`$ because $``$ is an abelian Lie group. Morever $`L_{}:\times G_{n,p}G_{n,p}`$ is a hamiltonian action and a momentum mapping for $`L_{}`$ is given by:
$$m:G_{n,p}L()^{}$$
$$\sigma \mu (\sigma )_{|L()}$$
where $`L_{}`$ is the Lie algebra of $``$. We arrive at the following theorem : 2.1 Theorem . Let $`(G_{n,p},d((H,N)^+)`$ defined as above. Then
1. $`m^1(H)`$ is a closed subgroup of $`G_{n,p}`$ and $`m^1(H)`$
2. The canonical exact sequence of Lie groups
$`(2)`$
$$\{ฯต\}m^1(H)m^1(H)/\{ฯต\}$$
is split. It is also an exact sequence of affine Lie groups.
3. The reduced symplectic Lie group $`m^1(H)/`$ is isomorphic to :
$$\begin{array}{cc}G_{np,p}(๐)& sin>p\\ G_{p\mathrm{1,1}}(๐)& sin=p>1\end{array}$$
4. In the principal bundle
$`(3)`$
$$m^1(H)\stackrel{i}{}G_{n,p}(๐)\stackrel{m}{}\mathrm{\Theta }$$
where $`\mathrm{\Theta }`$ is the set of matrices of rank $`p`$ in $`M_{n,p}^{}M_{p,n}`$, the fiber is an affine Lie subgroup and $`m`$ is affine relative to the usual affine structure of $`\mathrm{\Theta }๐^{np}`$.
We need the following lemma for which the proof is obvious. 2.2 Lemma . The mapping $`m`$ is a surjective submersion.
Furthermore,
$`(4)`$
$$m((X,T))=H.T^1\text{ where}(X,T)\text{is \hspace{1em}in }G_{n,p}(๐)$$
Proof. The formula (4) is a direct consequence of the definition of $`m`$ and the Lemma 1.2. Thus it is clear that $`m`$ is a surjective submersion on the set $`\mathrm{\Theta }`$ of matrices of rank $`p`$ of $`M_{n,p}(๐)`$. Proof of the theorem 2.1.
Formula (4) implies that $`m^1(H)=\{(X,T)G_{n,p}(๐)/HT^1=H\}`$ is a (closed) subgroup of $`G_{n,p}(๐)`$ which contains $``$. Morever the factor group $`m^1(H)/`$ can be identified with $`\{(0,T)G_{n,p}/HT^1=H\}`$. Thus (2) is an split sequence of Lie groups. On the other hand, since $``$ is commutative and $`\omega ^+`$ is exact, it follows that $`L()L()^{}`$ and a straightforward shows that $`L(m^1(H))=L()^{}`$.
Morever from the formula
$$\omega (xy,z)=\omega (y,[x,z])$$
defining the left symmetric product in $`๐ข_{n,p}(๐)`$ (see formula (1)) it turns out that $`L()^{}`$ is a left symmetric subalgebra of $`๐ข_{n,p}(๐)`$ and that $`L()`$ is a two-side ideal of $`L()^{}`$. Consequently (2) is an exact sequence of affine Lie groups i.e $``$, $`m^1(H)`$, $`m^1(H)/`$ are affine Lie groups and the applications of formula (2) are affine.
On the other hand a list of the elements of the matrix group $`\{(0,T)G_{n,p}(๐)/HT^1=H\}`$ allows to observe that the group $`m^1(H)/`$ is isomorphic to the group $`G_{np,p}(๐)`$ if $`n>p`$ and isomorphic to $`G_{n\mathrm{1,1}}(๐)`$ if $`n=p`$. This proves statement 3.
Now , we need to show that the manifold $`\mathrm{\Theta }`$ is endowed with an affine structure which makes the momentum mapping affine.
Let $`F`$ be the subbundle of $`TG_{n,p}(๐)`$ tangent to the $`L_{}`$-orbit and $`F^{}`$ its symplectic orthogonal. Denote by $``$ and $`^{}`$ the associated foliations respectively. Since $``$ is normal in $`G_{n,p}(๐)`$, the foliation $`^{}`$ is defined either by the left invariant form on $`G_{n,p}(๐)`$ given by $`\eta _j^{^{}}:=i(e_j^+)\omega ^+`$ where $`(e_i)`$โs form a basis of $`L()`$ and $`i`$ denote the interior product, or by the closed forms (thus exact) $`\eta _j:=i(e_j^{})\omega ^+`$. Obviously the $`\eta _j`$ are basic for the fibration. the forms $`\overline{\eta }_j`$ which are the projections of $`\eta _j`$ by $`m`$ define a local parallelism on $`\mathrm{\Theta }`$. This parallelism is global and commutative because the $`\overline{\eta }_j`$ are exact. Hence $`m`$ is affine. $`\mathrm{}`$
Let $`\alpha _1`$ and $`\alpha _2`$ be the linear forms on $`๐ข_{n,p}(๐)`$ defined respectively by $`\alpha _1(x,u)=tr(H.x)`$ and $`\alpha _2(x,u)=tr(N.u)`$. We get:
2.3. Theorem. $`ker(d\alpha _1)`$ and $`ker(d\alpha _2)`$ are supplementary symplectic Lie subalgebras of $`(๐ข_{n,p}(๐),d\alpha )`$. So they determine two tranverse symplectic foliations left invariant on $`G_{n,p}(๐)`$
Proof. Obviously $`ker(d\alpha _i)`$ is a Lie subalgebra of $`๐ข_{n,p}(๐)`$. In addition the subspaces $`ker(d\alpha _i)`$, $`i=\mathrm{1,2}`$ are in direct sum because ker$`(d\alpha _1)`$$``$ker$`(d\alpha _2)`$$`=\{0\}`$.
If we take $`\alpha =(H_0,N_0)`$ we directly observe that we have dim(ker$`(\delta \alpha _1))`$$`=p^2(k1)k`$ and dim(ker$`(\delta \alpha _2))`$$`=2p^2k`$. We extends these results about the dimension for all $`\alpha ๐ข_{n,p}^{}(๐)`$ with open coadjoint orbit. consequently $`ker(\delta \alpha _1)`$ and $`ker(\delta \alpha _2)`$ are symplectic Lie subalgebras of $`(๐ข_{n,p}(๐),d\alpha )`$. $`\mathrm{}`$ Let $`C(N_0)`$ be the subalgebra of $`gl(๐^n)`$ given by :
$$\left(\begin{array}{cccc}A_0& & & \\ A_1& \mathrm{}& & \\ \mathrm{}& \mathrm{}& \mathrm{}& \\ A_{k1}& \mathrm{}& A_1& A_0\end{array}\right)$$
2.4. Scholie. The Lie algebra $`ker(d\alpha _1)`$ is isomorphic to the Lie algebra $`๐ข_{np,p}(๐)`$ if $`n>p`$ and $`๐ข_{p\mathrm{1,1}}(๐)`$ if $`n=p`$ while $`ker(d\alpha _2)`$ is isomorphic to the semi-direct $`M_{n,p}(๐)C(N_0)`$ if $`n>p`$ and to the semi-direct product $`M_{p,p}(๐)C(N_0)`$ if $`n=p`$.
The following result is the infinitesimal version of Theorem 2.1. It gives a more precise and complete statement of Theorem 2.1. 2.5 Proposition. With the notations of Therem 2.1., if $`I=L()`$, the canonical sequence of vectoriel spaces
$$0II^{}I^{}/I0$$
is a split exact sequence of Lie algebras. It is also an exact sequence of left symmetric algebras.
Furthermore the Lie algebra $`๐ข_{n,p}(๐)`$ decomposes as a direct sum of Lie subalgebras $`I^{}`$ and $`C(N_0)`$ . 2.6 Theorem. The symplectic Lie group $`(G_{n,p},d\alpha ^+)`$ is endowed with two transversal left invariant symplectic foliations whose leaves are affine submanifolds of $`G_{n,p}(๐)`$.
The following assertion specifies the number of open orbits in $`๐ข_{n,p}^{}(๐)`$ as well as the left invariant symplectic structures on $`G_{n,p}(๐)`$. 2.7 Theorem.
a. There exist two open orbits of the coadjoint representation of $`G_{n,p}(๐)`$ if $`๐=`$ and only one if $`๐=`$
b. Up isomorphism there is only one left invariant symplectic structure on $`G_{n,p}`$ i.e if $`\omega `$ and $`\omega ^{^{}}`$ are two left invariant symplectic forms on $`G_{n,p}(๐)`$, then there exists an automorphism $`\phi `$ of Lie algebra $`๐ข_{n,p}(๐)`$ such that :
$$\omega _\epsilon (.,.)=\omega _\epsilon ^{^{}}(\phi .,\phi .).$$
The following lemmas set up the main steps of the demonstration of Theorem 2.7. This lemma allow to count the $`Ad_{G_{n,p}}^{}`$-orbits 2.8 Lemma. If $`Orb_{(H,M)}`$ is the coadjoint orbit of $`(H,M)๐ข_{n,p}^{}M_{p,n}(๐)\times gl(n)`$ then $`Orb_{(H,M)}`$ has an element $`(H_0^{^{}},M_0^{^{}})๐ข_{n,p}^{}`$ such that $`H_0^{^{}}=(0,\mathrm{},A_p)M_{p,n}`$.
Morever, if $`Orb_{(H,M)}`$ is open, then $`H_0^{^{}}`$ can be taken as $`H_0^{^{}}=(0,\mathrm{},I_p)=H_0`$. Proof. It suffices to remark that there exists $`UGL_o(๐^n)`$ such that $`HU^1=H_0`$ , that is obvious if we look at $`H`$ as the matrix of a mapping from $`๐^n`$ into $`๐^p`$ and $`U^1`$ as the matrix of a change of basis in $`๐^n`$. Then the last formula of lemma 1.2 show the first assertion.
Now suppose that $`Orb_{(H,M)}=Orb_{(H_0^{^{}},M_0^{^{}})}`$ is a open orbit. This implies that we have
$$(x,u)๐ข_{n,p};ad_{(x,u)}^{}(\alpha )=0(x,u)=0$$
In particular
$$ad_{(x\mathrm{,0})}^{}(H_0^{^{}},M^{^{}})=0x=0$$
In other words
$`()`$
$$tr(xH_0u)=0,ugl(n)x=0$$
with $`H_0=(0,A_p)`$
But a straight calculation shows that (\**) implies $`A_p`$ is invertible.
Finally there exists $`UGL_o(๐^n)`$ such that $`H_0^{^{}}U^1=H_0`$, but this relation fix only the last diagonal block of $`U`$ to the value $`A^1`$. Therefore if $`k2`$ we able to take an other diagonal block egals to det$`A^1.Id_{๐^p}`$ and completing the diagonale by the $`1`$, to construct a such matrix $`U`$ in $`SL(๐^n)`$. 2.9 Lemma. An open orbit $`Orb_{(H_0,M)}`$ contains only one element $`(H_0,M^{^{}})`$, where the block decomposition of $`M^{^{}}`$ can be written :
$$M^{^{}}=\left(\begin{array}{cc}M_1& 0\\ H_1& 0\end{array}\right)$$
with $`(H_1,M_1)๐ข_{np,p}^{}`$. Proof. Because
$$Ad_{(X,Id)}^{}(H_0,M)=(H_0,M+X.H)$$
it is clear that there is only one $`XM_{p,n}`$ such that $`M^{^{}}=M+MH_0`$. 2.10 Lemma. The linear forms $`(H_0,M^{^{}})`$ and $`(H_o,P^{^{}})`$ on $`๐ข_{n,p}`$ with $`M^{^{}}=(H_1,M_1)`$ and $`P^{^{}}=(K_1,P_1)`$, defined as in Lemma 2.9, belong to the same $`Ad_{G_{n,p}}^{}`$-orbit if and only if $`(H_1,M_1)`$ and $`(K_1,P_1)`$ are in the same $`Ad_{G_{np,p}}^{}`$-orbit. Proof. $`(H_0,M^{^{}})`$ and $`(H_o,P^{^{}})`$ are in the same orbit if and only if $`UGL_o(๐^n)`$ such that $`UM^{^{}}U^1=P^{^{}}`$ and $`H_0U^1=H_o`$.
The second condition implies
$$U=\left(\begin{array}{cc}U^{^{}}& 0\\ X_1& Id\end{array}\right),withU_1GL_0(๐^{np})$$
On other hand $`UM^{^{}}U^1=P^{^{}}`$ in $`gl(๐^n)`$ is equivalent to $`Ad_{(X_1,U_1)}^{}(H_1,M_1)=(K_1,P_1)`$ in $`๐ข_{np,p}(๐)`$ $`\mathrm{}`$
Proof of the theorem 2.7.
Following the previous Lemmas we have
cardinal$`\left\{openAd_{G_{n,p}}^{}orbit\right\}`$ = cardinal$`\left\{openAd_{G_{p,p}}^{}orbit\right\}`$
On the other hand, using similar arguments as those of the lemmas we can show that
cardinal$`\left\{openAd_{G_{p,p}}^{}orbit\right\}`$ = cardinal$`\left\{openAd_{G_{\mathrm{1,1}}}^{}orbit\right\}`$ =
$$\{\begin{array}{cc}2\hfill & \text{if}๐=\hfill \\ 1\hfill & \text{if}๐=\hfill \end{array}$$
If $`๐=`$ there is only one symplectic structure on $`G_{n,p}`$.
In the case $`๐=`$, it follows from the lemmas that every $`Ad_{G_{n,p}}^{}`$-open orbit contains an element $`(H_0,N_0)`$ where
$$N_0=\left(\begin{array}{ccccc}0& & & & \\ A_p& \mathrm{}& & & \\ & I_p& \mathrm{}& & \\ & & \mathrm{}& \mathrm{}& \\ & & & I_p& 0\end{array}\right)\text{ with }A_p\text{ invertible}$$
Nevertheless two matrices as $`N_0`$ are conjugate by an element $`P`$ of $`GL(n)`$.
However the mapping
$$๐ข_{n,p}^{}๐ข_{n,p}^{},(g,M)(P_{}^{t}{}_{}{}^{1}g,P^1MP)$$
is dual of the mapping
$$\theta :๐ข_{n,p}๐ข_{n,p}(x,N)(P^1x,P^1NP)$$
and these later is an automorphism of the Lie algebra $`๐ข_{n,p}`$.
Consequently if $`\omega _1^+`$ , $`\omega _2^+`$ are two left invariant symplectic forms on $`G_{n,p}`$, we have:
$$\omega _1^+=\theta ^{}(\omega _2^+).\mathrm{}$$
Remark 2. Notice that the only one left invariant affine structure on $`G_{n,p}`$ is given by $`(H_0,N_0)`$.
The following result is an important consequence of the previous study. 2.11 Proposition. The identity component of $`G_{n,p}`$ is diffeomorphic to an open $`Ad_{G_{n,p}}^{}`$-orbit. Consequently $`G_{n,p}`$ and $`\left\{\alpha ๐ข_{n,p}^{};\alpha \text{ is Poisson-regular}\right\}`$ are diffeomorphic. Proof.(By induction) We must prove that the orbital mapping in $`(H_0,N_0)`$
$$(G_{n,p})_o๐ข_{n,p}^{},(X,U)(H_0U^1,UN_0U^1+XH_0U^1)$$
has a trivial isotropy.
The result is obvious for $`G_{\mathrm{1,1}}`$ and this implies that it is also true for $`G_{p,p}`$ ( thanks to a sequence double extension ).
Consider the case $`G_{n,p}`$ with $`n=kp`$, $`k2`$. The equality
$$(H_0U^1,UN_0U^1+XH_0U^1)=(H_0,N_0)$$
implies that $`UN_0U^1`$ and $`N_0`$ have the same p block-type ( in particular their last columns are zero). Hence $`X=0`$.
On the other hand, $`U`$ induces an element of $`G_{np,p}`$ belonging to the $`Ad_{G_{np,p}}^{}`$-isotropy subgroup in $`N_0`$. Then, if the result is true in $`G_{p,p}`$ it is also true in $`G_{2p,p}`$ and $`G_{kp,p}`$$`p2`$.$`\mathrm{}`$ 3. Left invariant Lagrangian foliations and Hess connection .
3.1. Letโs specify a little bit the Lie algebra isomorphisms indicated by the previous proposition. The isomorphism between Rad$`(d\alpha _1)`$ and $`๐ข_{np,p}(๐)`$ is determined by the choice of a supplementary subspace of $`๐ฆ(H_0)=\{XM_{n,p}(๐),H_0.X=0\}`$ in $`M_{n,p}(๐)`$.
Denote by $`X_0`$, the element of $`M_{n,p}(๐)`$ formed by one column of zero-blocks except the last block which is $`id_{๐^p}`$. Then the mapping $`Rad(d\alpha _1)=\{(0,u)๐ข_{n,p}(๐),H_0.u=0\}๐ข_{np,p}(๐)`$ such that $`(0,u)(U.u_0,\pi _0(u))`$ defines such a isomorphism, where $`\pi _0(u)`$ denotes the matrix of the linear map given by $`u`$ restricted to $`๐ฆ(H_0)`$. Furthermore the image of the reduced symplectic form is the $`2`$-coboundary associated to $`(H_1,N_1)๐ข_{np,p}^{}(๐)`$ where $`H_1`$ and $`N_1`$ are defined by $`tr(N_0u)=tr(H_1.u.X_0)+tr(N_1\pi _0(u))`$ . We remark that $`(H_1,N_1)`$ has the same block type as $`(H_0,N_0)`$, so we can repeat the process of symplectic reduction in the same conditions. One obtains a decomposition of the space $`๐ข_{n,p}(๐)`$ as a direct sum of Lie subalgebras:
$$๐ข_{n,p}(๐)=๐ฆ_{p1}\mathrm{}๐ฆ_0C(N_0)\mathrm{}C(N_{p1})$$
where $`๐ฆ_i`$ comes from the totally isotropic ideal and $`๐ข_{nip,p}(๐)=\{K_i\times ๐ข_{n(i+1)p,p}(๐)\}C(N_i)`$ for $`0ip1`$ (orthogonal direct sum). Using the canonical embedding of $`๐ข_{n,p}(๐)`$ in $`gl(๐^{n+p})`$, the subspace $`L=๐ฆ_{p1}\mathrm{}๐ฆ_0`$ is identified to the subalgebra of strictly upper triangular $`(p+1)\times (p+1)`$-matrices by $`p\times p`$-blocks and $`L^{^{}}=C(N_0)\mathrm{}C(N_{p1})`$ to the subalgebra of lower triangular $`(p\times p`$)-matrices by $`(p\times p`$)-blocks. As the $`๐ฆ_i`$ and $`C(N_i)`$ appear as the totally isotropic subspaces at every step of the successive reductions the subalgebras $`L`$ and $`L^{}`$ are lagrangian relative to $`d\alpha `$ where $`\alpha (H_0,N_0)`$ .
Denote by $`\mathrm{\Lambda }`$ and $`\mathrm{\Lambda }^{^{}}`$ the connected Lie subgroups of $`G_{n,p}(๐)`$ with Lie algebra $`L`$ and $`L^{}`$ respectively. The natural left actions of $`\mathrm{\Lambda }`$ and $`\mathrm{\Lambda }^{^{}}`$ on $`(G_{n,p}(๐),d\alpha ^+)`$ being hamiltonnians, Theorem 3.1. of () allows us to assert that $`\mathrm{\Lambda }`$ and $`\mathrm{\Lambda }^{^{}}`$ are closed. So we have proved the following result:3.2 Theorem. The symplectic Lie group $`(G_{n,p},d\alpha ^+)`$ is endowed with two transversal left invariant lagrangian foliations with closed and affines leaves . Let $`(G_{n,p}(๐),\omega ^+)`$ be endowed with its affine structure $``$ defined by (1)
We recall that a connection on $`(G_{n,p}(๐),\omega ^+)`$ is said to be symplectic if and only if $`\omega =0`$ where $`\omega :=\omega _\epsilon ^+`$, in other words:
$$_a(\omega (b,c))=\omega (_ab,c)+\omega (b,_ac)a=(x,u),b=(y,v),c=(z,r)๐ข_{n,p}(๐)$$
In that follows one identifies an element $`a`$ of $`๐ข_{n,p}(๐)`$ with $`(x,u)`$, where $`xM_{n,p}`$ and $`ugl(n)`$.
Let $``$ and $`^{^{}}`$ be two lagrangian subalgebras of $`๐ข_{n,p}(๐)`$ such that $`๐ข=^{^{}}`$. Then we can write $`a=a_1+a_2`$ where $`a_1=(x_1,u_1)`$ and $`a_2=(x_2,u_2)^{^{}}`$ . Then the left symmetric product on $`๐ข_{n,p}(๐)`$ is given by :
$`(i)`$
$$(x\mathrm{,0}).(y\mathrm{,0})=(l,f)$$
where $`lM_{n,p}`$ is formed from a column of blocks such that the last belongs to $`sl(p)`$ and $`fgl(n)`$ is an element whose last line of blocks is zero.
$`(ii)`$
$$(0,u).(0,v)=(l,vu)$$
where $`l`$ is defined as above.
$`(iii)`$
$$(x\mathrm{,0}).(0,v)=(l_1,f_1)$$
$`(iv)`$
$$(0,u).(y\mathrm{,0})=(l_2,f_2)$$
where $`l_1`$,$`l_2M_{n,p}`$, $`f_1`$,$`f_2gl(n)`$ are defined as in (i).
The following result is a consequence of the previous discussion 3.3. Corollary. There exists only one (torsion free ) left invariant symplectic connection $`\stackrel{~}{}`$ ( called the Hessโconnection ) such that
$$\stackrel{~}{}_a,\stackrel{~}{}_a^{}^{^{}}^{^{}}.$$
where $`a`$ et $`a^{}^{^{}}`$
This connection is defined by the products :
$$\stackrel{~}{}_{(x\mathrm{,0})^+}(y\mathrm{,0})^+=(l\mathrm{,0})^+,\stackrel{~}{}_{(0,u)^+}(0,v)^+=(0,vu)^+$$
$$\stackrel{~}{}_{(x\mathrm{,0})^+}(0,v)^+=0,\stackrel{~}{}_{(0,u)^+}(y\mathrm{,0})^+=(uy\mathrm{,0})^+$$
Given the kind of results that we have described above, it is natural to ask ourself in what sense the group $`G_{n,p}(๐)`$ is not the symplectic double extension describes in ().
The answer is clearly no if $`k>1`$; in fact to have $`G_{n,p}(๐)`$ be the symplectic double extension of $`G_{np,p}(๐)`$ in the sense of it is necessary that $`I^{}`$ be a Lie ideal of $`๐ข_{n,p}(๐)`$. Considering the proposition 2.3 involve the existence of a Lie ideal of $`gl(๐^n)`$ isomorphic to $`๐ข_{np,p}(๐)`$. 4. $`G_{n,p}(๐)`$ as symplectic double extension of $`G_{np,p}(๐)`$
We have shown in the previous paragraph that the techniquess of symplectic double extension developed in () do not apply to $`G_{n,p}(๐)`$ for $`k>1`$.
We reconsider the study of the canonical fibrations (2) and (3 ) to try to understand how work this example. We have observed in 2.5. that for the symplectic Lie algebra $`(๐ข_{n,p},d\alpha )`$ where $`\alpha (H_0,N_0)`$, a section of the canonical exact sequence
$$0II^{}I^{}/I0$$
is determined by the choice of an element $`X_0`$ in $`M_{n,p}(๐)\backslash ๐ฆ(H_0)`$ where $`๐ฆ(H_0)`$ is the kernel of $`H_0`$.
Conversely, we consider the reduced algebra $`(๐ข_{np,p},d\alpha ^{^{}})`$ with $`\alpha ^{^{}}(H_1,N_1)`$ where $`H_1=(0,\mathrm{}\mathrm{,0},I_p)`$ and
$$N_1=\left(\begin{array}{ccccc}0& & & & \\ I_p& \mathrm{}& & & \\ & I_p& \mathrm{}& & \\ & & \mathrm{}& \mathrm{}& \\ & & & I_p& 0\end{array}\right)$$
Let $`i`$ be the canonical inclusion from $`M_{np,p}`$ to $`M_{n,p}`$ obtained by putting zero in the $`np+1,\mathrm{},n`$ rows, and let $`ZM_{n,p}\backslash M_{np,p}`$ have rank $`p`$ ( e.g. $`Z=H_0)`$).
Denote by $`r:M_{np,p}M_{n,p}`$ the linear mapping given by
$$ri=id_{M_{np,p}}\text{ and }r(Z)=0$$
and let $`HM_{n,p}^{}`$ satisfy
$`(5)`$
$$Hi=0\text{ and }H.Z=Id_p$$
We consider the regular representation :
$`(6)`$
$$\eta :๐ข_{np,p}gl(n),(x,u)(iur+H.i(x))$$
From $`\eta `$ we deduce a representation of Lie groups:
$`(7)`$
$$\rho :G_{np,p}GL(n),(X,U)(iUr+i(X).H)$$
verifying,
$$\rho _{,ฯต}=\eta $$
Morever consider the inclusion :
$$R:M_{n,p}^{}\times gl(np)gl(n)$$
given by
$`(8)`$
$$R(H,N)=iNr+Z.(Hr)$$
Then we can state the following result. 4.1 Theorem. Consider the symplectic Lie group $`(G_{np,p},d(H_1,N_1)^+)`$ . If $`H`$ is the linear form given by (4), and $`N=R(H_1,N_1)`$ with $`R`$ given by (7), then $`(G_{n,p},d(H,N)^+)`$ is a symplectic Lie group such that $`(G_{np,p},d(H_1,N_1)^+)`$ is the reduced symplectic Lie group as described in theorem 2.1. Proof. By definition of $`H`$ we have
$$m^1(H)=M_{n,p}\times \rho (G_{np,p})$$
It remains to prouve that $`d(H_1,N_1)^+)`$ is the reduced symplectic form of $`d(H,N)^+)_{|m^1(H)}`$. Consider the decomposition $`M_{n,p}(๐)=i(M_{np,p}(๐))๐.Z`$
On the other hand a straight verification proves that
$`tr(H_1x)+tr(N_1u)=tr(R(H_1,N_1)(\eta (x,u))`$ for all $`(x,u)`$ of $`๐ข_{np,p}(๐)`$. Finally as $`N=R(H_1,N_1)`$ and $`d((H\mathrm{,0})^+)`$ vanish on $`m^1(H)`$ it follows that $`d((H_1,N_1))^+)`$ is the reduced form .This ends the proof. $`\mathrm{}`$ Remark 4. A similar argument proves that for $`p2`$ : $`G_{p\mathrm{,1}}=Aff(๐^p)`$ is a symplectic reduction of $`G_{p,p}`$ and $`G_{p,p}`$ is a symplectic double extension of $`G_{p\mathrm{,1}}`$ 4.2 Study of the fibration :
$`(3)`$
$$m^1(H)\stackrel{i}{}GA_{n,p}(๐)\stackrel{m}{}\mathrm{\Theta }$$
According to Theorem 1.2. we can consider the case where $`\alpha (H_0,N_0)`$. Let $`C=C(N_0)GL(๐^n)`$ or, if we prefer $`C`$ is the Lie subgroup of $`GL(๐^n)`$ formed by matrices of type
$$\left(\begin{array}{cccc}A_0& & & \\ A_1& \mathrm{}& & \\ \mathrm{}& \mathrm{}& \mathrm{}& \\ A_{k1}& \mathrm{}& A_1& A_0\end{array}\right)$$
with $`A_0`$ is invertible.
According to Lemma 2.2. $`m(C)=\{(A_{k1},\mathrm{},A_1,A_0):A_0\text{ invertible }\}`$, denote by $`V_0`$ this open set of $`\mathrm{\Theta }`$.
Let be $`V_\gamma =V_0.\sigma _\gamma `$ where $`\gamma =(i_1,\mathrm{},i_p)`$ and $`\sigma _\gamma `$ is the element of $`GL(๐^n)`$ which realizes the permutation of the $`k`$ last columns with the columns indexed by $`\gamma `$ in $`M_{p,n}(๐)`$. The sets $`V_\gamma `$ are clearly the open set of $`\mathcal{\Theta }`$ and $`m`$ defines a diffeomorphism of $`C.\sigma _\gamma `$ on $`V_\gamma `$. Denote $`S_\gamma `$ the embedding of $`V_\gamma `$ in $`GL(๐^n)`$ such that $`S_\gamma (V_\gamma )=C.\sigma _\gamma `$ and $`mS_\gamma =id_{U_\gamma }`$ for all multi-indices $`\gamma `$.
We have the following result: 4.3 Lemma. The $`V_\gamma `$ as defined above , with $`\gamma =(i_1,\mathrm{},i_p)`$ for $`1i_1<i_2<\mathrm{}<i_pn`$, form an open trivializing cover for the fibration (2). Proof.
The elements $`V_\gamma `$ recover $`\mathrm{\Theta }`$ since $`\mathrm{\Theta }`$ is formed by matrices of rank $`p`$ in $`M_{p,n}(๐)`$ whose first $`p`$ columns are independent in $`๐^p`$. By contrast $`V_0`$ is formed of matrices for which the last $`p`$ colums are independent.
To show that the $`V_\gamma `$ are trivialising for the fibration (3) amounts to proving that we have $`m^1(V_\gamma )=S_\gamma (V_\gamma ).m^1(H_0)`$. For all $`\sigma `$ of $`G_{n,p}(๐)`$ we have the identity $`m^1(m(\sigma ))=\sigma .m^1(H_0)`$. Indeed if $`\sigma ^{^{}}m^1(m(\sigma ))`$ we have by definition of $`m`$ the formula,
$$Ad_\sigma ^{}(H\mathrm{,0});(X\mathrm{,0})=Ad_\sigma ^{^{}}^{}(H\mathrm{,0});(X\mathrm{,0})\text{ for all }X\text{ in }M_{n,p}(๐)$$
or what amounts to the same thing
$$Ad_{\sigma ^1\sigma ^{^{}}}^{}(H\mathrm{,0});(X\mathrm{,0})=(H\mathrm{,0});(X\mathrm{,0})\text{ for all }X\text{ in }M_{n,p}(๐)$$
This last relation means that $`\sigma ^1\sigma ^{^{}}m^1(H_0)`$ by the identity below and the fact that the $`V_\gamma `$ are open and trivialising.The trivialisations are then given by the map $`\varphi _\gamma :m^1(V_\gamma )V_\gamma \times m^1(H_0);S_\gamma (\alpha ).\sigma (\alpha ,\sigma )`$ for all multi-indices $`\gamma =(i_1,\mathrm{},i_p)`$ with $`i_1<i_2<K<i_pn`$. The cocycle defining the fibration (2) is given by
$`(9)`$
$$\mathrm{\Gamma }_{\gamma _1\gamma _2}:V_{\gamma _1}V_{\gamma _2}m^1(H_0);\alpha S_{\gamma _1}^1(\alpha ).S_{\gamma _2}(\alpha )$$
Thus we have proved the following result. 4.4 Proposition. The manifold $`G_{n,p}(๐)`$ is diffeomorphic to $`V_\gamma \times m^1(H_0)/`$ where $`(\alpha ,\sigma )(\beta ,\tau )`$ if and only if $`\alpha =\beta `$ and $`\sigma =\mathrm{\Gamma }_{\gamma _1\gamma _2}(\tau )`$ for $`(\alpha ,\sigma )`$ in $`V_{\gamma _1}\times m^1(H_0)`$ and $`(\beta ,\tau )`$ in $`V_{\gamma _2}\times m^1(H_0)`$. The cocycle $`\mathrm{\Gamma }_{\gamma _1\gamma _2}`$ is defined by (9).
Acknowledgements. The third author was partially supported by NSF grant DMS-0204100 and The Department of Mathematics of Penn state University. He also wishes to thank The Department of Romances Languages and the Department of Mathematics at The University of North Carolina at Chapel Hill for their kind welcome during his stay. Thanks to Jim Stasheff for editing this paper and thanks to Patrick Eberlein for his helpful discussions and suggestions.
$`\begin{array}{cc}JeanMichelDARDIE\hfill & AlbertoMEDINA\hfill \\ Universit\text{รฉ}deMontpellier2\hfill & Universit\text{รฉ}deMontpellier2\hfill \\ D\text{รฉ}partementdeMath\text{รฉ}matiques\hfill & D\text{รฉ}partementdeMath\text{รฉ}matiques\hfill \\ 34095Montpelliercedex5\hfill & 34095Montpelliercedex5\hfill \\ UMR\mathrm{\hspace{0.17em}5149}duCNRS\hfill & UMR\mathrm{\hspace{0.17em}5149}duCNRS\hfill \\ dardie\mathrm{@}darboux.math.univmontp2.fr\hfill & medina\mathrm{@}darboux.math.univmontp2.fr\hfill \end{array}`$
$$\begin{array}{c}HasseneSIBY\hfill \\ Universit\text{รฉ}deMontpellier2\hfill \\ D\text{รฉ}partementdeMath\text{รฉ}matiques\hfill \\ 34095Montpelliercedex5\hfill \\ UMR\mathrm{\hspace{0.17em}5149}duCNRS\hfill \\ siby\mathrm{@}darboux.math.univmontp2.fr\hfill \end{array}$$
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# Invariant dโHermite des jacobiennes de graphes pondรฉrรฉs
## Introduction
Soit $`n`$ un entier non nul. Lโรฉtude de la densitรฉ des rรฉseaux dโun espace euclidien $`(^n,<.,.>)`$ est un sujet classique et prend la forme suivante : รฉtant donnรฉ un rรฉseau $`\mathrm{\Lambda }^n`$, nous dรฉfinissons le dรฉterminant de $`\mathrm{\Lambda }`$, notรฉ $`det(\mathrm{\Lambda })`$, comme le carrรฉ du volume euclidien du domaine fondamental du rรฉseau et sa norme minimale par la formule
$$||\mathrm{\Lambda }||=\mathrm{min}\{<\lambda ,\lambda >\lambda \mathrm{\Lambda }\{0\}\}.$$
Lโinvariant dโHermite du rรฉseau est la quantitรฉ
$$\mu (\mathrm{\Lambda })=\frac{\mathrm{\Lambda }}{\sqrt[n]{det(\mathrm{\Lambda })}}$$
et code la densitรฉ du rรฉseau. La densitรฉ maximale en dimension $`n`$ correspond ร la constante dยดHermite :
$$\gamma _n=sup\{\mu (\mathrm{\Lambda })\mathrm{\Lambda }\text{ rรฉseau de }^n\}.$$
(1)
Cette quantitรฉ est bien dรฉfinie et vรฉrifie lโencadrement suivant (voir ) :
$$\frac{n}{2\pi e}\gamma _n\frac{1.744n}{2\pi e}.$$
(2)
Dans la dรฉfinition (1), on peut considรฉrer la borne supรฉrieure non plus sur tous les rรฉseaux mais sur un sous-ensemble de rรฉseaux de $`^n`$. P. Buser et P. Sarnak ont รฉtudiรฉ dans la borne supรฉrieure des invariants dโHermite des rรฉseaux symplectiques et ont montrรฉ quโelle vรฉrifiait lโinรฉgalitรฉ infรฉrieure dans la formule (2). Ils ont รฉgalement montrรฉ le rรฉsultat suivant : la borne supรฉrieure des invariants dโHermite sur lโensemble des rรฉseaux formรฉ des jacobiennes de surfaces de Riemann de genre $`g`$, que lโon notera $`\eta _{2g}`$, vรฉrifie
$$c\mathrm{ln}g\eta _{2g}\frac{3}{\pi }\mathrm{ln}(4g+3),$$
(3)
$`c`$ est une constante positive.
Le but de cet article est de prouver un rรฉsultat analogue pour la jacobienne des graphes pondรฉrรฉs. Rappelons tout dโabord quelques dรฉfinitions. Un graphe $`\mathrm{\Gamma }=(V,E)`$ est un complexe simplicial de dimension 1. Cโest la donnรฉe dโune paire dโensembles $`(V,E)`$, oรน $`V`$ dรฉsigne les sommets et $`E`$ les arรชtes. La valence dโun sommet est le nombre dโarรชtes incidentes ร ce sommet et un graphe sera dit $`k`$-rรฉgulier, pour $`k^{}`$, si la valence de chacun de ses sommets est constante รฉgale ร $`k`$. Dans ce qui suit, les graphes seront supposรฉs connexes, finis et leurs sommets de valence au moins $`2`$. Un graphe pondรฉrรฉ est une paire $`(\mathrm{\Gamma },w)`$$`\mathrm{\Gamma }=(V,E)`$ est un graphe et $`w`$ est une fonction poids sur les arรชtes $`w:E_+`$. Nous dirons dโun graphe pondรฉrรฉ $`(\mathrm{\Gamma },w)`$ quโil est combinatoire si sa fonction poids est constante รฉgale ร $`1`$, et nous le noterons simplement $`\mathrm{\Gamma }`$. Le type dโhomotopie dโun graphe donnรฉ $`\mathrm{\Gamma }=(V,E)`$ est caractรฉrisรฉ par le nombre $`b_1(\mathrm{\Gamma })`$ de cycles indรฉpendants, appelรฉ premier nombre de Betti ou nombre cyclomatique. On a la formule $`b_1(\mathrm{\Gamma })=|E||V|+1`$$`|X|`$ dรฉsigne le cardinal dโun ensemble fini $`X`$. Notons quโร premier nombre de Betti fixรฉ, les graphes considรฉrรฉs sont en nombre fini ร homรฉomorphisme prรจs.
Etant donnรฉ un graphe pondรฉrรฉ $`(\mathrm{\Gamma },w)`$ de premier nombre de Betti $`b1`$, nous choisissons pour chaque arรชte une orientation arbitraire et nous noterons $`๐ผ=\{e_i\}_{i=1}^k`$ lโensemble de ces arรชtes orientรฉes. Soit
$$๐(\mathrm{\Gamma },)=\{\underset{i=1}{\overset{k}{}}a_i.e_ia_i\text{ pour }i=1,\mathrm{},k\}$$
lโespace vectoriel engendrรฉ par les arรชtes orientรฉes. Cet espace coรฏncide avec lโespace des chaรฎnes simpliciales du complexe simplicial $`\mathrm{\Gamma }`$. Il est muni du produit scalaire naturel $`<e_i,e_j>_w=w(e_i)\delta _{ij}`$ pour $`1ijk`$, oรน $`\delta _{ij}`$ dรฉsigne le symbole de Kronecker (voir , p. 191). Lโhomologie de $`\mathrm{\Gamma }`$ de dimension $`1`$ ร coefficients rรฉels $`H_1(\mathrm{\Gamma },)`$ est plongรฉe dans $`๐(\mathrm{\Gamma },)`$ comme un sous-espace vectoriel de dimension $`b`$ et on note encore $`<.,.>_w`$ la restriction du produit scalaire ร ce sous-espace. Lโhomologie de $`\mathrm{\Gamma }`$ de dimension $`1`$ ร coefficients entiers $`H_1(\mathrm{\Gamma },)`$, en lโabsence de torsion dans ce cadre unidimensionnel, constitue un rรฉseau du sous-espace $`H_1(\mathrm{\Gamma },)`$ (comparer avec ).
On note $`\mathrm{\Lambda }(\mathrm{\Gamma },w)=H_1(\mathrm{\Gamma },)(H_1(\mathrm{\Gamma },),<.,.>_w)`$ le rรฉseau de dimension $`b`$ ainsi dรฉfini et on lโappelle jacobienne du graphe pondรฉrรฉ $`(\mathrm{\Gamma },w)`$ : il est clair que lโinvariant dโHermite de ce rรฉseau ne dรฉpend pas des orientations dโarรชtes initialement choisies. Posons
$$\rho _b=sup\{\mu (\mathrm{\Lambda }(\mathrm{\Gamma },w))(\mathrm{\Gamma },w)\text{ graphe pondรฉrรฉ de premier nombre de Betti }b\}.$$
Notre rรฉsultat principal sโรฉnonce alors de la maniรจre suivante.
Thรฉorรจme Pour $`b`$ suffisamment grand,
$$\frac{1}{6e}\mathrm{log}_2b\rho _b4\mathrm{log}_2(\frac{8}{3}b),$$
(4)
$`\mathrm{log}_2`$ dรฉsigne le logarithme en base $`2`$.
Nous pouvons, sous certaines restrictions, amรฉliorer ces deux inรฉgalitรฉs.
1) Pour une infinitรฉ de valeurs $`\{b_m\}_m`$, il existe un graphe combinatoire $`3`$-rรฉguliers $`G_m`$ de premier nombre de Betti $`b_m`$ (construit dans ) pour lequel (voir inรฉgalitรฉ (9))
$$\mu (\mathrm{\Lambda }(G_m))\frac{4}{9e}\mathrm{log}_2b_m,$$
2) Tout graphe combinatoire $`\mathrm{\Gamma }`$, dont la valence en chaque sommet est au moins $`3`$, vรฉrifie (voir inรฉgalitรฉ (10))
$$\mu (\mathrm{\Lambda }(\mathrm{\Gamma }))2\mathrm{log}_2b.$$
La suite de cet article est consacrรฉ aux dรฉmonstrations de ces rรฉsultats. Nous allons montrer que lโรฉtude de lโinvariant dโHermite de la jacobienne dโun graphe pondรฉrรฉ est รฉquivalente ร lโรฉtude du problรจme combinatoire suivant : รฉtant donnรฉ un graphe, borner la complexitรฉ de ce graphe par sa systole (voir la section $`1`$ pour les dรฉfinitions de ces quantitรฉs). Dans la premiรจre partie de ce papier, nous donnons la dรฉmonstration de lโinรฉgalitรฉ supรฉrieure dans la formule (4). Dans la seconde partie, nous majorons pour tout graphe sa complexitรฉ par son volume unidimensionnel, et obtenons ainsi ร lโaide de graphes systoliquement รฉconomiques - graphes dont le rapport volume sur systole est suffisamment petit - lโestimรฉe infรฉrieure annoncรฉe pour $`\rho _b`$. Nous dรฉmontrons pour finir les amรฉliorations 1) et 2).
## 1 Complexitรฉ et systole dโun graphe
Nous allons montrer la proposition suivante, qui dรฉmontrera la majoration de $`\rho _b`$ annoncรฉe dans le thรฉorรจme par un calcul dโรฉquivalent รฉlรฉmentaire.
###### Proposition 1
Pour tout graphe pondรฉrรฉ $`(\mathrm{\Gamma },w)`$ de premier nombre de Betti $`b2`$,
$$\mu (\mathrm{\Lambda }(\mathrm{\Gamma },w))4\left(\underset{k=2}{\overset{b}{}}\mathrm{log}_2(\frac{8}{3}k)\right)^{1/b}.$$
(5)
Dรฉmonstration. Commenรงons par rรฉduire le problรจme. On peut construire facilement un graphe $`3`$-rรฉgulier $`\mathrm{\Gamma }^{}=(V^{},E^{})`$ et une fonction poids $`w^{}`$ sur $`\mathrm{\Gamma }^{}`$ tels que $`(\mathrm{\Gamma },w)`$ soit obtenu ร partir de $`(\mathrm{\Gamma }^{},w^{})`$ en contractant les arรชtes de poids nul en un point. Notons que $`\mu (\mathrm{\Lambda }(\mathrm{\Gamma },w))=\mu (\mathrm{\Lambda }(\mathrm{\Gamma }^{},w^{}))`$. Comme lโapplication
$`\mu _\mathrm{\Gamma }^{}:^{|E^{}|}`$ $``$ $`_+`$
$`w`$ $``$ $`\mu (\mathrm{\Lambda }(\mathrm{\Gamma }^{},w))`$
est continue et invariante par composition avec les dilatations, on en dรฉduit que pour tout $`ฯต>0`$, il existe une fonction poids $`w_ฯต`$ sur $`\mathrm{\Gamma }^{}`$ telle que :
\- Pour toute arรชte $`eE^{},w_ฯต(e)`$,
\- $`|\mu (\mathrm{\Lambda }(\mathrm{\Gamma },w))\mu (\mathrm{\Lambda }(\mathrm{\Gamma }^{},w_ฯต))|<ฯต`$.
Il suffit donc de dรฉmontrer le rรฉsultat annoncรฉ pour un graphe combinatoire ($`w=1`$) dont les sommets sont de valence $`2`$ ou $`3`$.
Etant donnรฉ un graphe combinatoire $`\mathrm{\Gamma }`$, nous pouvons exprimer lโinvariant dโHermite de sa jacobienne en fonction des quantitรฉs suivantes :
\- La systole, qui dans le cas dโun graphe pondรฉrรฉ $`(\mathrm{\Gamma },w)`$ est dรฉfinie comme la plus petite longueur dโun circuit simple de $`\mathrm{\Gamma }`$ et est notรฉe $`\text{sys}(\mathrm{\Gamma },w)`$. Dans le cas dโun graphe $`\mathrm{\Gamma }`$ combinatoire, nous noterons simplement $`\text{sys}(\mathrm{\Gamma })`$ cette quantitรฉ.
\- La complexitรฉ, qui est dรฉfinie comme le nombre dโarbres maximaux du graphe $`\mathrm{\Gamma }`$ (voir ) et est notรฉe $`\kappa (\mathrm{\Gamma })`$.
Dโaprรจs , รฉtant donnรฉ un graphe combinatoire $`\mathrm{\Gamma }`$,
$`\mathrm{\Lambda }(\mathrm{\Gamma })=\text{sys}(\mathrm{\Gamma })`$ et $`det(\mathrm{\Lambda }(\mathrm{\Gamma }))=\kappa (\mathrm{\Gamma })`$.
On en dรฉduit lโรฉgalitรฉ
$$\mu (\mathrm{\Lambda }(\mathrm{\Gamma }))=\frac{\text{sys}(\mathrm{\Gamma })}{\sqrt[b]{\kappa (\mathrm{\Gamma })}}.$$
Supposons pour la suite de cette dรฉmonstration que $`\mathrm{\Gamma }`$ est un graphe combinatoire dont les sommets sont de valence $`2`$ ou $`3`$. Notre majoration revient donc ร estimer infรฉrieurement la complexitรฉ dโun tel graphe par sa systole. Nous obtenons le rรฉsultat suivant, qui implique lโestimรฉe (5) :
###### Lemme 1
$$\kappa (\mathrm{\Gamma })\frac{\text{sys}(\mathrm{\Gamma })^b}{4^b_{k=2}^b\mathrm{log}_2(\frac{8}{3}k)}.$$
(6)
Dรฉmonstration du lemme. Tout dโabord, supposons que $`b=2`$. Comme les sommets de $`\mathrm{\Gamma }`$ sont de valence $`2`$ ou $`3`$, on a deux classes dโhomรฉomorphismes possibles pour $`\mathrm{\Gamma }`$ : la classe $`8_1`$ et la classe $`8_2`$ (voir figure 1).
Il est clair que $`\kappa (8_1)(\text{sys}(8_1))^2`$ et $`\kappa (8_2)(\text{sys}(8_2)/2)^2`$, dโoรน (6) dans ce cas.
On suppose maintenant $`b>2`$. Si
$$\left[\frac{\text{sys}(\mathrm{\Gamma })}{2\mathrm{log}_2(\frac{8}{3}b)}\right]=0,$$
$`[n]`$ dรฉsigne la partie entiรจre dโun entier $`n`$, le rรฉsultat est immรฉdiat comme la complexitรฉ dโun graphe est un entier non nul. Supposons donc
$$\left[\frac{\text{sys}(\mathrm{\Gamma })}{2\mathrm{log}_2(\frac{8}{3}b)}\right]1.$$
On note $`\stackrel{~}{\mathrm{\Gamma }}`$ le graphe dรฉfini ร partir de $`\mathrm{\Gamma }`$ de la maniรจre suivante. Etant donnรฉ un sommet $`v`$ de valence $`2`$, on considรจre le graphe obtenu en supprimant le sommet $`v`$ et les deux arรชtes $`e_1`$ et $`e_2`$ incidentes ร ce sommet, et en ajoutant une nouvelle arรชte reliant les deux sommets de $`e_1`$ et $`e_2`$ restants (voir figure $`2`$).
On rรฉpรจte lโopรฉration pour tout sommet de valence $`2`$ et on obtient ainsi le graphe $`\stackrel{~}{\mathrm{\Gamma }}`$. Cโest un graphe $`3`$-rรฉgulier et on note $`f:\mathrm{\Gamma }\stackrel{~}{\mathrm{\Gamma }}`$ lโapplication topologique naturelle qui envoie uniformรฉment une suite maximale dโarรชtes adjacentes de $`\mathrm{\Gamma }`$ dont les sommets intermรฉdiaires sont de valence $`2`$ sur lโarรชte correspondante de $`\stackrel{~}{\mathrm{\Gamma }}`$. Soit $`\stackrel{~}{\gamma }`$ une courbe rรฉalisant sa systole et $`\gamma `$ son image inverse par $`f`$. La systole est estimรฉe supรฉrieurement pour les graphes $`3`$-rรฉguliers de premier nombre de Betti $`b`$ de la maniรจre suivante (voir ) :
$$\text{sys}(\stackrel{~}{\mathrm{\Gamma }})[2\mathrm{log}_2(\frac{2}{3}b)+3].$$
Donc il existe une arรชte $`e\stackrel{~}{\gamma }\stackrel{~}{\mathrm{\Gamma }}`$ telle que, si $`C=f^1(e)`$, on ait
$$\text{long}(C)\left[\frac{\text{sys}(\mathrm{\Gamma })}{2\mathrm{log}_2(\frac{8}{3}b)}\right].$$
En effet, sinon, $`\text{long}(\gamma )<\text{sys}(\mathrm{\Gamma })`$, dโoรน une contradiction.
Soit $`\mathrm{\Gamma }C`$ le complรฉmentaire de la suite dโarรชtes adjacentes $`C`$ dans $`\mathrm{\Gamma }`$ : cโest un sous-graphe de $`\mathrm{\Gamma }`$ vรฉrifiant $`b_1(\mathrm{\Gamma }C)=b_1(\mathrm{\Gamma })1`$. Tout arbre maximal $`T`$ du graphe $`\mathrm{\Gamma }C`$ fournit de maniรจre รฉvidente au moins $`[\text{sys}(\mathrm{\Gamma })/(2\mathrm{log}_2(\frac{8}{3}b))]`$ arbres maximaux $`T^{}`$ de $`\mathrm{\Gamma }`$, et tous les arbres maximaux ainsi obtenus, lorsque lโon fait varier lโarbre initial $`T`$ parmi les arbres maximaux de $`\mathrm{\Gamma }C`$, sont deux ร deux distincts. On en dรฉduit :
$$\kappa (\mathrm{\Gamma })\kappa (\mathrm{\Gamma }C).\left[\frac{\text{sys}(\mathrm{\Gamma })}{2\mathrm{log}_2(\frac{8}{3}b)}\right]\kappa (\mathrm{\Gamma }C).\frac{\text{sys}(\mathrm{\Gamma })}{2(2\mathrm{log}_2(\frac{8}{3}b))}.$$
Comme $`\text{sys}(\mathrm{\Gamma }C)\text{sys}(\mathrm{\Gamma })`$, on obtient le rรฉsultat par rรฉcurrence. $`\mathrm{}`$
Remarque. Nous avons donc reliรฉ la complexitรฉ et la systole de tout graphe $`\mathrm{\Gamma }`$ de premier nombre de Betti $`b2`$ de la maniรจre suivante :
$$\sqrt[b]{\kappa (\mathrm{\Gamma })}\frac{\text{sys}(\mathrm{\Gamma })}{4\left(_{k=2}^b(\mathrm{log}_2(\frac{8}{3}k))\right)^{1/b}}\frac{\text{sys}(\mathrm{\Gamma })}{4\mathrm{log}_2(\frac{4}{3}b)}.$$
De lโinรฉgalitรฉ (8), nous dรฉduisons que lโordre de grandeur dans cette estimรฉe infรฉrieure est optimale.
## 2 Complexitรฉ et volume dโun graphe
Nous allons dans cette section dรฉfinir la notion de constante systolique pour un graphe pondรฉrรฉ, et ainsi pouvoir introduire une famille de graphes systoliquement รฉconomiques. Cette famille nous permet, ร lโaide dโune majoration de la complexitรฉ dโun graphe par son volume, de dรฉmontrer lโinรฉgalitรฉ infรฉrieure annoncรฉe dans la formule (4).
Etant donnรฉ un graphe pondรฉrรฉ $`(\mathrm{\Gamma },w)`$ de premier nombre de Betti $`b2`$, le problรจme systolique peut รชtre formulรฉ comme suit. On dรฉfinit la constante systolique de $`\mathrm{\Gamma }`$ par
$$\sigma (\mathrm{\Gamma })=\underset{w}{inf}\frac{\text{Vol}(\mathrm{\Gamma },w)}{\text{sys}(\mathrm{\Gamma },w)}$$
$`\text{Vol}(\mathrm{\Gamma },w)`$ dรฉsigne le $`1`$-volume de $`\mathrm{\Gamma }`$ (la somme des poids de ses arรชtes) et oรน lโinfimum est pris sur lโensemble des fonctions poids de $`\mathrm{\Gamma }`$. Nous avons, dโaprรจs ,
$$\sigma (\mathrm{\Gamma })\frac{3}{2}\frac{b1}{\mathrm{log}_2(b1)+\mathrm{log}_2\mathrm{log}_2(b1)+4}.$$
(7)
Dโautre part, comme il a รฉtรฉ dรฉmontrรฉ dans , il existe pour chaque $`b2`$ un graphe combinatoire de premier nombre de Betti $`b`$ que nous noterons $`\mathrm{\Gamma }_b`$ et que nous appelerons systoliquement รฉconomique vรฉrifiant asymptotiquement
$$\frac{\text{Vol}(\mathrm{\Gamma }_b)}{\text{sys}(\mathrm{\Gamma }_b)}6\frac{b}{\mathrm{log}_2b},$$
$`\text{Vol}(\mathrm{\Gamma }_b)`$ est le $`1`$-volume de $`\mathrm{\Gamma }_b`$ pour la fonction poids constante รฉgale ร $`1`$.
Nous pouvons minorer lโinvariant dโHermite de la jacobienne dโun graphe pondรฉrรฉ par sa constante systolique comme suit :
###### Proposition 2
Pour tout graphe pondรฉrรฉ $`(\mathrm{\Gamma },w)`$ de premier nombre de Betti $`b1`$, on a
$$\mu (\mathrm{\Lambda }(\mathrm{\Gamma },w))\frac{\sqrt[b]{b!}}{\sigma (\mathrm{\Gamma },w)}.$$
Dรฉmonstration. De mรชme que dans le dรฉbut de la section 1, nous nous ramenons ร un graphe combinatoire. On estime supรฉrieurement la complexitรฉ de ce graphe par son volume de la maniรจre suivante. Tout arbre maximal $`T`$ de $`\mathrm{\Gamma }`$ est entiรจrement dรฉterminรฉ par les $`b`$ arรชtes de $`\mathrm{\Gamma }T`$, dโoรน
$$\kappa (\mathrm{\Gamma })C_{\text{Vol}(\mathrm{\Gamma })}^b\frac{\text{Vol}(\mathrm{\Gamma })^b}{b!}.$$
On en dรฉduit immรฉdiatement le rรฉsultat. $`\mathrm{}`$
Les graphes $`\mathrm{\Gamma }_b`$ vรฉrifient donc
$$\mu (\mathrm{\Lambda }(\mathrm{\Gamma }_b))\frac{1}{6}\frac{\sqrt[b]{b!}}{b}\mathrm{log}_2b,$$
dโoรน la minoration annoncรฉe en (4) par application de la formule de Stirling :
$$\mu (\mathrm{\Lambda }(\mathrm{\Gamma }_b))\frac{1}{6e}\mathrm{log}_2b.$$
(8)
Nous expliquons maintenant les amรฉliorations 1) et 2) annoncรฉes dans lโintroduction.
1) Nous pouvons, en utilisant les familles de graphes ร grand tour de taille mises en รฉvidence dans , amรฉliorer pour certaines valeurs de $`b`$ lโestimรฉe (8). G.A. Margulis a construit, pour une famille infinie $`\{b_m\}_m`$ de valeurs, une famille $`\{G_m\}_m`$ de graphes $`3`$-rรฉguliers de premier nombre de Betti $`\{b_m\}_m`$ pour lesquels
$$\text{sys}(G_m)\frac{4}{3}\mathrm{log}_2b_m.$$
Nous en dรฉduisons
$$\mu (\mathrm{\Lambda }(G_m))\frac{4}{9e}\mathrm{log}_2b_m.$$
(9)
2) Si lโon se restreint ร certains graphes, nous pouvons estimer supรฉrieurement lโinvariant dโHermite de la jacobienne par la constante systolique :
###### Proposition 3
Pour tout graphe combinatoire $`\mathrm{\Gamma }`$ dont chaque sommet est de valence au moins trois et de premier nombre de Betti $`b2`$,
$$\mu (\mathrm{\Lambda }(\mathrm{\Gamma }))\frac{3(b1)}{\sigma (\mathrm{\Gamma })}.$$
Dรฉmonstration. Comme la valence en chaque sommet du graphe est au moins $`3`$, on a lโinรฉgalitรฉ
$$\text{Vol}(\mathrm{\Gamma })3(b1),$$
donc
$$\mu (\mathrm{\Lambda }(\mathrm{\Gamma }))=\frac{\text{sys}(\mathrm{\Gamma })}{\sqrt[b]{\kappa (\mathrm{\Gamma })}}\text{sys}(\mathrm{\Gamma })\text{Vol}(\mathrm{\Gamma })\frac{\text{sys}(\mathrm{\Gamma })}{\text{Vol}(\mathrm{\Gamma })}\frac{3(b1)}{\sigma (\mathrm{\Gamma })}.$$
$`\mathrm{}`$
On en dรฉduit une amรฉlioration de lโinรฉgalitรฉ (5) pour cette classe de graphes, en vertu de lโestimรฉe systolique (7) : pour tout graphe $`\mathrm{\Gamma }`$ dont chaque sommet est de valence au moins trois et de premier nombre de Betti $`b2`$,
$$\mu (\mathrm{\Lambda }(\mathrm{\Gamma }))2(\mathrm{log}_2(b1)+\mathrm{log}_2\mathrm{log}_2(b1)+4)2\mathrm{log}_2b.$$
(10)
Remerciements. Je remercie Roland Bacher de mโavoir suggรฉrer cette รฉtude, ainsi que Ivan Babenko pour les nombreuses discussions que nous avons partagรฉ.
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# Spin precession in the Schwarzschild spacetime: circular orbits
## 1 Introduction
The equations of motion for a spinning test particle in a given gravitational background were deduced by Mathisson and Papapetrou and read
$`{\displaystyle \frac{DP^\mu }{\mathrm{d}\tau _U}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}R^\mu {}_{\nu \alpha \beta }{}^{}U_{}^{\nu }S^{\alpha \beta }F^{(\mathrm{sc})}{}_{}{}^{\mu },`$ (1.1)
$`{\displaystyle \frac{DS^{\mu \nu }}{\mathrm{d}\tau _U}}`$ $`=`$ $`P^\mu U^\nu P^\nu U^\mu ,`$ (1.2)
where $`P^\mu `$ is the total four-momentum of the particle, $`S^{\mu \nu }`$ is the antisymmetric spin tensor, and $`U`$ is the 4-velocity of the timelike โcenter of mass lineโ used to make the multipole reduction. In order to have a closed set of equations, Eqs. (1.1) and (1.2) must be completed by adding supplementary conditions (SC) whose standard choices in the literature are the
* Corinaldesi-Papapetrou conditions (CP): $`S^{t\nu }=0`$, where the index $`\nu `$ corresponds to a coordinate component and $`t`$ is a timelike slicing coordinate,
* Pirani conditions (P): $`S^{\mu \nu }U_\nu =0`$,
* Tulczyjew conditions (T): $`S^{\mu \nu }P_\nu =0`$.
Only solutions of the combined equations for which both $`U`$ and $`P`$ are timelike vectors are considered, in order to have a meaningful interpretation describing a spinning test particle with nonzero rest mass and physical momentum.
Not much is known about actual solutions of these equations in explicit spacetimes which satisfy the Einstein equations. In a previous article , we considered the simplest special case of a spinning test particle moving uniformly along a circular orbit in the static spherically symmetric Schwarzschild spacetime, but because these equations are still complicated, we looked for solutions with constant frame components of the spin tensor in the natural symmetry adapted static frame, i.e., coinciding with a static tensor field along the path. Such a static spin tensor is a very strong restriction on the solutions of these equations of motion, leading to special solutions in which the spin vector is perpendicular to the plane of the orbit, and contributes to an adjustment in the acceleration of the orbit.
Here we consider the slightly less restrictive case where the spin components are not constant, but the motion is still circular. However, in this case it is clear that if the spin tensor has time-dependent components, its feedback into the acceleration of the test particle path will break the static symmetry of that path unless the spin precession is very closely tied to the natural Frenet-Serret rotational properties of the path itself. Indeed we find that only the Pirani supplementary conditions permit such specialized solutions since they allow the spin tensor to be described completely by a spatial spin vector in the local rest space of the path itself. By locking spin vector precession to the Frenet-Serret rotational velocity of the path, solutions are found with a spin vector Fermi-Walker transported along an accelerated center of mass world line. The remaining choices for the supplementary conditions have no natural relationship to the Frenet-Serret properties of the particle path and do not admit such specialized solutions.
With the assumption of circular motion, one can solve the equations of motion explicitly up to constants of integration. By a process of elimination, one can express them entirely in terms of the spin components and particle mass as a constant coefficient linear system of first and second order differential equations. By systematic solving and backsubstitution, one gets decoupled linear second order constant coefficient equations for certain spin components, which are easily solved to yield exponential or sinusoidal or quadratic solutions as functions of the proper time, from which the remaining variables may be calculated. Imposing the choice of supplementary conditions then puts constraints on the constants of integration or leads to inconsistencies. The details of the decoupling and solution of the equations of motion are left to the Appendix, leaving the imposition of the supplementary conditions to the main text.
## 2 Circular orbits in the Schwarzschild spacetime
Consider the case of the Schwarzschild spacetime, with the metric written in standard coordinates
$$\mathrm{d}s^2=\left(1\frac{2M}{r}\right)\mathrm{d}t^2+\left(1\frac{2M}{r}\right)^1\mathrm{d}r^2+r^2(\mathrm{d}\theta ^2+\mathrm{sin}^2\theta \mathrm{d}\varphi ^2),$$
(2.1)
and introduce the usual orthonormal frame adapted to the static observers following the time lines
$$e_{\widehat{t}}=(12M/r)^{1/2}_t,e_{\widehat{r}}=(12M/r)^{1/2}_r,e_{\widehat{\theta }}=\frac{1}{r}_\theta ,e_{\widehat{\varphi }}=\frac{1}{r\mathrm{sin}\theta }_\varphi ,$$
(2.2)
with dual frame
$$\omega ^{\widehat{t}}=(12M/r)^{1/2}\mathrm{d}t,\omega ^{\widehat{r}}=(12M/r)^{1/2}\mathrm{d}r,\omega ^{\widehat{\theta }}=r\mathrm{d}\theta ,\omega ^{\widehat{\varphi }}=r\mathrm{sin}\theta \mathrm{d}\varphi ,$$
(2.3)
where $`\{_t,_r,_\theta ,_\varphi \}`$ and $`\{\mathrm{d}t,\mathrm{d}r,\mathrm{d}\theta ,\mathrm{d}\varphi \}`$ are the coordinate basis and its dual, respectively.
In order to investigate the simplest special solutions of the combined equations of motion, we explore the consequences of assuming that the test particle 4-velocity $`U`$ corresponds to a timelike constant speed circular orbit confined to the equatorial plane $`\theta =\pi /2`$. Then it must have the form
$$U=\mathrm{\Gamma }[_t+\zeta _\varphi ]=\gamma [e_{\widehat{t}}+\nu e_{\widehat{\varphi }}],\gamma =(1\nu ^2)^{1/2},$$
(2.4)
where $`\zeta `$ is the angular velocity with respect to infinity, $`\nu `$ is the azimuthal velocity as seen by the static observers, $`\gamma `$ is the associated gamma factor, and $`\mathrm{\Gamma }`$ is a normalization factor which assures that $`UU=1`$. These are related by
$$\zeta =(g_{tt}/g_{\varphi \varphi })^{1/2}\nu ,\mathrm{\Gamma }=\left(g_{tt}\zeta ^2g_{\varphi \varphi }\right)^{1/2}=(g_{tt})^{1/2}\gamma ,$$
(2.5)
so that $`\zeta \mathrm{\Gamma }=\gamma \nu /(g_{\varphi \varphi })^{1/2}`$, which reduces to $`\gamma \nu /r`$ in the equatorial plane.
Here $`\zeta `$ and therefore $`\nu `$ are assumed to be constant along the world line. We limit our analysis to the equatorial plane $`\theta =\pi /2`$; as a convention, the physical (orthonormal) component along $`_\theta `$ which is perpendicular to the equatorial plane will be referred to as โalong the positive $`z`$-axisโ and will be indicated by the index $`\widehat{z}`$ when convenient: $`e_{\widehat{z}}=e_{\widehat{\theta }}`$. Note both $`\theta =\pi /2`$ and $`r=r_0`$ are constants along any given circular orbit, and that the azimuthal coordinate along the orbit depends on the coordinate time $`t`$ or proper time $`\tau `$ along that orbit according to
$$\varphi \varphi _0=\zeta t=\mathrm{\Omega }_U\tau _U,\mathrm{\Omega }_U=\gamma \nu /r,$$
(2.6)
defining the corresponding coordinate and proper time orbital angular velocities $`\zeta `$ and $`\mathrm{\Omega }_U`$. These determine the rotation of the spherical frame with respect to a nonrotating frame at infinity.
Among all circular orbits the timelike circular geodesics merit special attention, whether co-rotating $`(\zeta _+)`$ or counter-rotating $`(\zeta _{})`$ with respect to increasing values of the azimuthal coordinate $`\varphi `$ (counter-clockwise motion). Their time coordinate angular velocities $`\zeta _\pm \pm \zeta _K=\pm (M/r^3)^{1/2}`$, which are identical with the Newtonian Keplerian values, lead to the expressions
$$U_\pm =\gamma _K[e_{\widehat{t}}\pm \nu _Ke_{\widehat{\varphi }}],\nu _K=\left[\frac{M}{r2M}\right]^{1/2},\gamma _K=\left[\frac{r2M}{r3M}\right]^{1/2},$$
(2.7)
where the timelike condition $`\nu _K<1`$ is satisfied if $`r>3M`$. At $`r=3M`$ these circular geodesics go null.
It is convenient to introduce the Lie relative curvature of each orbit
$$k_{(\mathrm{lie})}=_{\widehat{r}}\mathrm{ln}\sqrt{g_{\varphi \varphi }}=\frac{1}{r}\left(1\frac{2M}{r}\right)^{1/2}=\frac{\zeta _K}{\nu _K},$$
(2.8)
and a Frenet-Serret intrinsic frame along $`U`$ , defined by
$$E_0=U,E_1=e_{\widehat{r}},E_2=\gamma [\nu e_{\widehat{t}}+e_{\widehat{\varphi }}],E_3=e_{\widehat{z}}$$
(2.9)
satisfying the following system of evolution equations along the constant radial acceleration orbit
$$\frac{DU}{d\tau _U}a(U)=\kappa E_1,\frac{DE_1}{d\tau _U}=\kappa U+\tau _1E_2,\frac{DE_2}{d\tau _U}=\tau _1E_1,\frac{DE_3}{d\tau _U}=0,$$
(2.10)
where in this case
$`\kappa `$ $`=`$ $`k_{(\mathrm{lie})}\gamma ^2[\nu ^2\nu _K^2]={\displaystyle \frac{\gamma ^2(\nu ^2\nu _K^2)}{\nu _K}}\zeta _K,`$
$`\tau _1`$ $`=`$ $`{\displaystyle \frac{1}{2\gamma ^2}}{\displaystyle \frac{d\kappa }{d\nu }}=k_{(\mathrm{lie})}{\displaystyle \frac{\gamma ^2}{\gamma _K^2}}\nu ={\displaystyle \frac{\gamma ^2\nu }{\gamma _K^2\nu _K}}\zeta _K.`$ (2.11)
The projection of the spin tensor into the local rest space of the static observers defines the spin vector by spatial duality
$$S^\beta =\frac{1}{2}\eta _\alpha {}_{}{}^{\beta \gamma \delta }(e_{\widehat{t}})_{}^{\alpha }S_{\gamma \delta },$$
(2.12)
where $`\eta _{\alpha \beta \gamma \delta }=\sqrt{g}ฯต_{\alpha \beta \gamma \delta }`$ is the unit volume 4-form constructed from the Levi-Civita alternating symbol $`ฯต_{\alpha \beta \gamma \delta }`$ ($`ฯต_{\widehat{t}\widehat{r}\widehat{\theta }\widehat{\varphi }}=1`$), leading to the correspondence
$$(S^{\widehat{r}},S^{\widehat{\theta }}=S^{\widehat{z}},S^{\widehat{\varphi }})=(S_{\widehat{\theta }\widehat{\varphi }},S_{\widehat{r}\widehat{\varphi }},S_{\widehat{r}\widehat{\theta }}).$$
(2.13)
For the CP supplementary conditions only these components of the spin tensor remain nonzero, while in the remaining cases the other nonzero components are determined from these through the corresponding orthogonality condition. The total spin scalar is also useful
$$s^2=\frac{1}{2}S_{\mu \nu }S^{\mu \nu }=S_{\widehat{r}\widehat{t}}^2S_{\widehat{\theta }\widehat{t}}^2S_{\widehat{\varphi }\widehat{t}}^2+S_{\widehat{r}\widehat{\theta }}^2+S_{\widehat{r}\widehat{\varphi }}^2+S_{\widehat{\theta }\widehat{\varphi }}^2,$$
(2.14)
and in general is not constant along the trajectory of the spinning particle. In the Schwarzschild field the total spin must be small enough compared to the mass of the test particle and of the black hole $`|s|/(mM)1`$ for the approximation of the Mathisson-Papapetrou model to be valid. This inequality follows from requiring that the characteristic length scale $`|s|/m`$ associated with the particleโs internal structure be small compared to the natural length scale $`M`$ associated with the background field in order that the particle backreaction can be neglected, i.e., that the description of a test particle on a background field make sense .
## 3 Solving the equations of motion: preliminary steps
Consider first the evolution equation for the spin tensor (1.2). By contracting both sides of Eq. (1.2) with $`U_\nu `$, one obtains the following expression for the total 4-momentum
$$P^\mu =(UP)U^\mu U_\nu \frac{DS^{\mu \nu }}{\mathrm{d}\tau _U}mU^\mu +P_s^\mu ,$$
(3.1)
which then defines the particleโs mass $`m`$, which a priori does not have to be constant along the orbit, while $`P_s^\mu =U_\alpha DS^{\alpha \mu }/\mathrm{d}\tau _U`$ is the part of the 4-momentum orthogonal to $`U`$. Finally let $`U_p`$ denote the timelike unit vector associated with the total 4-momentum $`P=PU_p\mu U_p`$.
Backsubstituting this representation Eq. (3.1) of the momentum into the spin evolution Eq. (1.2) expressed in the static observer frame leads to
$`0`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}S_{\widehat{r}\widehat{\varphi }}}{\mathrm{d}\tau _U}}\nu {\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{r}}}{\mathrm{d}\tau _U}}+\gamma {\displaystyle \frac{\zeta _K}{\nu _K}}(\nu ^2\nu _K^2)S_{\widehat{t}\widehat{\varphi }},`$ (3.2)
$`0`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}S_{\widehat{\theta }\widehat{\varphi }}}{\mathrm{d}\tau _U}}\nu {\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{\theta }}}{\mathrm{d}\tau _U}}{\displaystyle \frac{\gamma \nu }{\gamma _K^2}}{\displaystyle \frac{\zeta _K}{\nu _K}}S_{\widehat{r}\widehat{\theta }},`$ (3.3)
$`0`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}S_{\widehat{r}\widehat{\theta }}}{\mathrm{d}\tau _U}}\gamma \nu _K\zeta _KS_{\widehat{t}\widehat{\theta }}+\gamma \nu {\displaystyle \frac{\zeta _K}{\nu _K}}S_{\widehat{\theta }\widehat{\varphi }}.`$ (3.4)
From (3.1), using the definition of $`P_s`$ and equations (3.2)โ(3.4), it follows that the total 4-momentum $`P`$ can be written in the form
$`P`$ $`=`$ $`\gamma (m+\nu m_s)e_{\widehat{t}}+{\displaystyle \frac{1}{\gamma }}\left[{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{r}}}{\mathrm{d}\tau _U}}\gamma \nu {\displaystyle \frac{\zeta _K}{\nu _K}}S_{\widehat{t}\widehat{\varphi }}\right]e_{\widehat{r}}+{\displaystyle \frac{1}{\gamma }}\left[{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{\theta }}}{\mathrm{d}\tau _U}}\gamma \nu _K\zeta _KS_{\widehat{r}\widehat{\theta }}\right]e_{\widehat{\theta }}`$
$`+\gamma (m\nu +m_s)e_{\widehat{\varphi }}`$
$`=`$ $`mU+m_sE_{\widehat{\varphi }}+{\displaystyle \frac{1}{\gamma }}\left[{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{r}}}{\mathrm{d}\tau _U}}\gamma \nu {\displaystyle \frac{\zeta _K}{\nu _K}}S_{\widehat{t}\widehat{\varphi }}\right]e_{\widehat{r}}+{\displaystyle \frac{1}{\gamma }}\left[{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{\theta }}}{\mathrm{d}\tau _U}}\gamma \nu _K\zeta _KS_{\widehat{r}\widehat{\theta }}\right]e_{\widehat{\theta }},`$ (3.5)
with
$$m_s=\frac{\mathrm{d}S_{\widehat{t}\widehat{\varphi }}}{\mathrm{d}\tau _U}+\gamma \nu \frac{\zeta _K}{\nu _K}S_{\widehat{t}\widehat{r}}\gamma \nu _K\zeta _KS_{\widehat{r}\widehat{\varphi }}.$$
(3.6)
Next consider the equation of motion (1.1). The Riemann tensor spin-curvature-coupling force term is
$`F^{(\mathrm{sc})}=\gamma \zeta _K^2\left\{\nu S_{\widehat{t}\widehat{\varphi }}e_{\widehat{t}}+\left[2S_{\widehat{t}\widehat{r}}+\nu S_{\widehat{r}\widehat{\varphi }}\right]e_{\widehat{r}}\left[S_{\widehat{t}\widehat{\theta }}+2\nu S_{\widehat{\theta }\widehat{\varphi }}\right]e_{\widehat{\theta }}S_{\widehat{t}\widehat{\varphi }}e_{\widehat{\varphi }}\right\}.`$ (3.7)
Using (3.1), the balance condition which allows a circular orbit of this type to exist can be written as
$$ma(U)=F^{(\mathrm{so})}+F^{(\mathrm{sc})},$$
(3.8)
where $`a(U)`$ is the acceleration of the $`U`$-orbit and $`F^{(\mathrm{so})}DP_s/d\tau _U`$ defines the spin-orbit coupling force term, which arises from the variation of the spin along the orbit.
Taking (3) and (3.6) into account, Eq. (1.1) gives rise to the following set of ordinary differential equations
$`0`$ $`=`$ $`\nu {\displaystyle \frac{\mathrm{d}^2S_{\widehat{t}\widehat{\varphi }}}{\mathrm{d}\tau _U^2}}2\gamma \nu \nu _K\zeta _K{\displaystyle \frac{\mathrm{d}S_{\widehat{r}\widehat{\varphi }}}{\mathrm{d}\tau _U}}+\gamma {\displaystyle \frac{\zeta _K}{\nu _K}}(\nu ^2+\nu _K^2){\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{r}}}{\mathrm{d}\tau _U}}`$
$`\gamma ^2\nu \zeta _K^2(\nu ^2\nu _K^2)S_{\widehat{t}\widehat{\varphi }}+{\displaystyle \frac{\mathrm{d}m}{\mathrm{d}\tau _U}},`$ (3.9)
$`0`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}^2S_{\widehat{t}\widehat{r}}}{\mathrm{d}\tau _U^2}}\nu {\displaystyle \frac{\mathrm{d}^2S_{\widehat{r}\widehat{\varphi }}}{\mathrm{d}\tau _U^2}}2{\displaystyle \frac{\gamma \nu }{\gamma _K^2}}{\displaystyle \frac{\zeta _K}{\nu _K}}{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{\varphi }}}{\mathrm{d}\tau _U}}+\nu \gamma ^2\zeta _K^2(\nu ^2\nu _K^2)S_{\widehat{r}\widehat{\varphi }}`$
$`\zeta _K^2\left[{\displaystyle \frac{\gamma ^2}{\gamma _K^2}}{\displaystyle \frac{\nu ^2}{\nu _K^2}}+2\right]S_{\widehat{t}\widehat{r}}m\gamma {\displaystyle \frac{\zeta _K}{\nu _K}}(\nu ^2\nu _K^2),`$ (3.10)
$`0`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}^2S_{\widehat{t}\widehat{\theta }}}{\mathrm{d}\tau _U^2}}\nu {\displaystyle \frac{\mathrm{d}^2S_{\widehat{\theta }\widehat{\varphi }}}{\mathrm{d}\tau _U^2}}+\gamma {\displaystyle \frac{\zeta _K}{\nu _K}}(\nu ^2\nu _K^2){\displaystyle \frac{\mathrm{d}S_{\widehat{r}\widehat{\theta }}}{\mathrm{d}\tau _U}}+2\nu \zeta _K^2S_{\widehat{\theta }\widehat{\varphi }}+\zeta _K^2S_{\widehat{t}\widehat{\theta }},`$ (3.11)
$`0`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}^2S_{\widehat{t}\widehat{\varphi }}}{\mathrm{d}\tau _U^2}}\gamma {\displaystyle \frac{\zeta _K}{\nu _K}}(\nu ^2+\nu _K^2){\displaystyle \frac{\mathrm{d}S_{\widehat{r}\widehat{\varphi }}}{\mathrm{d}\tau _U}}+2\gamma \nu {\displaystyle \frac{\zeta _K}{\nu _K}}{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{r}}}{\mathrm{d}\tau _U}}\zeta _K^2\left[{\displaystyle \frac{\gamma ^2}{\gamma _K^2}}{\displaystyle \frac{\nu ^2}{\nu _K^2}}1\right]S_{\widehat{t}\widehat{\varphi }}`$
$`+\nu {\displaystyle \frac{\mathrm{d}m}{\mathrm{d}\tau _U}}.`$ (3.12)
Note that there are two equations containing the second derivative of $`S_{\widehat{t}\widehat{\varphi }}`$; this is due to the presence of its first derivative in two different components of $`P`$ (more precisely, in $`P^{\widehat{t}}`$ and $`P^{\widehat{\varphi }}`$, see Eqs. (3) and (3.6)).
Once the system of constant coefficient linear differential equations (3.2)โ(3.4) and (3)โ(3) is solved for $`m`$ and the spin tensor components, one may then calculate $`P`$. The system must be decoupled, leading to functions which are either exponentials, sinusoidals, or at most quadratic functions of the proper time along the particle world line. The elimination method for decoupling the equations is crucially different depending on whether $`\nu `$ has the values 0 or $`\pm \nu _K`$ or none of these values, since one or the other or neither term drops out of the spin equations (3.2) and so must be considered separately. From the details of their derivations discussed in the Appendix, one sees why there are several zones approaching the horizon where the solutions change character.
## 4 Particles at rest: the $`\nu =0`$ case
When the particle is at rest, the solutions for the components of the spin tensor and the varying mass $`m`$ of the spinning particle are given by
1. $`2M<r<3M`$:
$`S_{\widehat{\theta }\widehat{\varphi }}`$ $`=`$ $`c_1,`$
$`S_{\widehat{t}\widehat{r}}`$ $`=`$ $`c_2e^{\omega _1\tau }+c_3e^{\omega _1\tau }+{\displaystyle \frac{\nu _K}{\zeta _K}}{\displaystyle \frac{c_m}{2+\nu _K^2}},`$
$`m`$ $`=`$ $`\nu _K\zeta _K[c_2e^{\omega _1\tau }+c_3e^{\omega _1\tau }]+{\displaystyle \frac{2c_m}{2+\nu _K^2}},`$
$`S_{\widehat{t}\widehat{\theta }}`$ $`=`$ $`c_4e^{\overline{\omega }_0\tau }+c_5e^{\overline{\omega }_0\tau },`$
$`S_{\widehat{t}\widehat{\varphi }}`$ $`=`$ $`c_6e^{\overline{\omega }_0\tau }+c_7e^{\overline{\omega }_0\tau },`$
$`S_{\widehat{r}\widehat{\theta }}`$ $`=`$ $`\gamma _K\nu _K\left[c_4e^{\overline{\omega }_0\tau }c_5e^{\overline{\omega }_0\tau }\right]+c_8,`$
$`S_{\widehat{r}\widehat{\varphi }}`$ $`=`$ $`\gamma _K\nu _K\left[c_6e^{\overline{\omega }_0\tau }c_7e^{\overline{\omega }_0\tau }\right]+c_9;`$ (4.1)
2. $`r=3M`$:
$`S_{\widehat{\theta }\widehat{\varphi }}`$ $`=`$ $`c_1,`$
$`S_{\widehat{t}\widehat{r}}`$ $`=`$ $`c_2e^{\tau /(3M)}+c_3e^{\tau /(3M)}+\sqrt{3}Mc_m,`$
$`m`$ $`=`$ $`{\displaystyle \frac{\sqrt{3}}{9M}}[c_2e^{\tau /(3M)}+c_3e^{\tau /(3M)}]+{\displaystyle \frac{2}{3}}c_m,`$
$`S_{\widehat{t}\widehat{\theta }}`$ $`=`$ $`c_4\tau +c_5,`$
$`S_{\widehat{t}\widehat{\varphi }}`$ $`=`$ $`c_6\tau +c_7,`$
$`S_{\widehat{r}\widehat{\theta }}`$ $`=`$ $`{\displaystyle \frac{\sqrt{3}}{9M}}\left[{\displaystyle \frac{c_4}{2}}\tau ^2+c_5\tau \right]+c_8,`$
$`S_{\widehat{r}\widehat{\varphi }}`$ $`=`$ $`{\displaystyle \frac{\sqrt{3}}{9M}}\left[{\displaystyle \frac{c_6}{2}}\tau ^2+c_7\tau \right]+c_9;`$ (4.2)
3. $`r>3M`$:
$`S_{\widehat{\theta }\widehat{\varphi }}`$ $`=`$ $`c_1,`$
$`S_{\widehat{t}\widehat{r}}`$ $`=`$ $`c_2e^{\omega _1\tau }+c_3e^{\omega _1\tau }+{\displaystyle \frac{\nu _K}{\zeta _K}}{\displaystyle \frac{c_m}{2+\nu _K^2}},`$
$`m`$ $`=`$ $`\nu _K\zeta _K[c_2e^{\omega _1\tau }+c_3e^{\omega _1\tau }]+{\displaystyle \frac{2c_m}{2+\nu _K^2}},`$
$`S_{\widehat{t}\widehat{\theta }}`$ $`=`$ $`c_4\mathrm{cos}\omega _0\tau +c_5\mathrm{sin}\omega _0\tau ,`$
$`S_{\widehat{t}\widehat{\varphi }}`$ $`=`$ $`c_6\mathrm{cos}\omega _0\tau +c_7\mathrm{sin}\omega _0\tau ,`$
$`S_{\widehat{r}\widehat{\theta }}`$ $`=`$ $`\gamma _K\nu _K\left[c_4\mathrm{sin}\omega _0\tau c_5\mathrm{cos}\omega _0\tau \right]+c_8,`$
$`S_{\widehat{r}\widehat{\varphi }}`$ $`=`$ $`\gamma _K\nu _K\left[c_6\mathrm{sin}\omega _0\tau c_7\mathrm{cos}\omega _0\tau \right]+c_9,`$ (4.3)
where $`c_m,c_1,\mathrm{},c_9`$ are integration constants and
$$\omega _0=i\overline{\omega }_0=\frac{\zeta _K}{\gamma _K}=\sqrt{\frac{M(r3M)}{r^3(r2M)}},\omega _1=\zeta _K(2+\nu _K^2)^{1/2}=\sqrt{\frac{M(2r3M)}{r^3(r2M)}}.$$
(4.4)
From Eq. (3) the total 4-momentum $`P`$ then has the value
$$P=me_{\widehat{t}}+\omega _1[c_2e^{\omega _1\tau }c_3^{\omega _1\tau }]e_{\widehat{r}}\frac{\zeta _K}{\nu _K}\left[S_{\widehat{r}\widehat{\theta }}\frac{c_8}{\gamma _K^2}\right]e_{\widehat{\theta }}\frac{\zeta _K}{\nu _K}\left[S_{\widehat{r}\widehat{\varphi }}\frac{c_9}{\gamma _K^2}\right]e_{\widehat{\varphi }}$$
(4.5)
in cases (i) and (iii), and
$$P=me_{\widehat{t}}+\frac{1}{3M}[c_2e^{\tau /(3M)}c_3^{\tau /(3M)}]e_{\widehat{r}}\left[\frac{\sqrt{3}}{9M}S_{\widehat{r}\widehat{\theta }}c_4\right]e_{\widehat{\theta }}\left[\frac{\sqrt{3}}{9M}S_{\widehat{r}\widehat{\varphi }}c_6\right]e_{\widehat{\varphi }}$$
(4.6)
in case (ii).
At this point the supplementary conditions impose constraints on the constants of integration which appear in the solution. For a particle at rest ($`\nu =0`$), the CP and P conditions coincide and imply that $`S_{\widehat{t}\widehat{a}}=0`$, namely
$$c_2=c_3=c_4=c_5=c_6=c_7=0,c_m=0,$$
(4.7)
leaving arbitrary values for $`c_1,c_8,c_9`$. As a consequence, $`m`$ should be $`0`$ as well, implying that $`P`$ should be spacelike and therefore physically inconsistent.
The T supplementary conditions when $`\nu =0`$ imply instead
$`0`$ $`=`$ $`S_{\widehat{t}\widehat{\theta }}{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{\theta }}}{\mathrm{d}\tau _U}}+S_{\widehat{t}\widehat{r}}{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{r}}}{\mathrm{d}\tau _U}}+S_{\widehat{t}\widehat{\varphi }}{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{\varphi }}}{\mathrm{d}\tau _U}}\nu _K\zeta _K[S_{\widehat{t}\widehat{\varphi }}S_{\widehat{r}\widehat{\varphi }}+S_{\widehat{t}\widehat{\theta }}S_{\widehat{r}\widehat{\theta }}],`$
$`0`$ $`=`$ $`S_{\widehat{r}\widehat{\theta }}{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{\theta }}}{\mathrm{d}\tau _U}}+S_{\widehat{r}\widehat{\varphi }}{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{\varphi }}}{\mathrm{d}\tau _U}}\nu _K\zeta _K(S_{\widehat{r}\widehat{\theta }}^2+S_{\widehat{r}\widehat{\varphi }}^2)mS_{\widehat{t}\widehat{r}},`$
$`0`$ $`=`$ $`S_{\widehat{\theta }\widehat{\varphi }}{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{\varphi }}}{\mathrm{d}\tau _U}}S_{\widehat{r}\widehat{\theta }}{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{r}}}{\mathrm{d}\tau _U}}\nu _K\zeta _KS_{\widehat{\theta }\widehat{\varphi }}S_{\widehat{r}\widehat{\varphi }}mS_{\widehat{t}\widehat{\theta }},`$
$`0`$ $`=`$ $`S_{\widehat{\theta }\widehat{\varphi }}{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{\theta }}}{\mathrm{d}\tau _U}}+S_{\widehat{r}\widehat{\varphi }}{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{r}}}{\mathrm{d}\tau _U}}\nu _K\zeta _KS_{\widehat{r}\widehat{\theta }}S_{\widehat{\theta }\widehat{\varphi }}+mS_{\widehat{t}\widehat{\varphi }}.`$ (4.8)
By substituting the solutions given by Eqs. (1)โ(3) into these equations, one finds that all the integration constants except $`c_1`$ must vanish. This in turn implies $`m=0`$, which again leads to a spacelike $`P`$. Thus a spinning particle with nonzero rest mass cannot remain at rest in the given gravitational field.
## 5 Geodesic motion: the $`\nu =\pm \nu _K`$ case
When the test particleโs center of mass moves along a geodesic (the orbit has zero acceleration $`a(U)=0`$) with azimuthal velocity $`\nu =\pm \nu _K`$, the spin-curvature and the spin-orbit forces balance each other (see Eq. (3.8)): $`F_{(\mathrm{so})}=F^{(\mathrm{sc})}`$. The solution of Eqs. (3)-(3) determines the spin which leads to this balancing. In the Schwarzschild spacetime, timelike circular geodesics only exist for $`r>3M`$. We consider separately the various cases:
1. $`3M<r<6M`$:
$`S_{\widehat{\theta }\widehat{\varphi }}`$ $`=`$ $`c_3\mathrm{cos}\omega _4\tau +c_4\mathrm{sin}\omega _4\tau ,`$
$`S_{\widehat{t}\widehat{\theta }}`$ $`=`$ $`c_5\mathrm{cos}\omega _3\tau +c_6\mathrm{sin}\omega _3\tau \pm {\displaystyle \frac{S_{\widehat{\theta }\widehat{\varphi }}}{\nu _K}},`$
$`S_{\widehat{r}\widehat{\theta }}`$ $`=`$ $`\gamma _K\nu _K\left[c_5\mathrm{sin}\omega _3\tau c_6\mathrm{cos}\omega _3\tau \right],`$
$`m`$ $`=`$ $`c_m,`$
$`S_{\widehat{t}\widehat{r}}`$ $`=`$ $`c_7e^{\overline{\omega }_2\tau }+c_8e^{\overline{\omega }_2\tau }\pm 2{\displaystyle \frac{\gamma _K}{\zeta _K}}{\displaystyle \frac{c_2}{43\gamma _K^2}},`$
$`S_{\widehat{r}\widehat{\varphi }}`$ $`=`$ $`\pm \nu _K[c_7e^{\overline{\omega }_2\tau }+c_8e^{\overline{\omega }_2\tau }]+c_1+2{\displaystyle \frac{\nu _K}{\zeta _K}}\gamma _K{\displaystyle \frac{c_2}{43\gamma _K^2}},`$
$`S_{\widehat{t}\widehat{\varphi }}`$ $`=`$ $`\pm {\displaystyle \frac{2}{\gamma _K}}(3\gamma _K^24)^{1/2}[c_8e^{\overline{\omega }_2\tau }c_7e^{\overline{\omega }_2\tau }]+c_93\gamma _K^2{\displaystyle \frac{c_2}{43\gamma _K^2}}\tau ;`$ (5.1)
2. $`r=6M`$:
$`S_{\widehat{\theta }\widehat{\varphi }}`$ $`=`$ $`c_3\mathrm{cos}{\displaystyle \frac{\sqrt{3}\tau }{18M}}+c_4\mathrm{sin}{\displaystyle \frac{\sqrt{3}\tau }{18M}},`$
$`S_{\widehat{t}\widehat{\theta }}`$ $`=`$ $`c_5\mathrm{cos}{\displaystyle \frac{\sqrt{6}\tau }{36M}}+c_6\mathrm{sin}{\displaystyle \frac{\sqrt{6}\tau }{36M}}\pm 2S_{\widehat{\theta }\widehat{\varphi }},`$
$`S_{\widehat{r}\widehat{\theta }}`$ $`=`$ $`{\displaystyle \frac{\sqrt{3}}{3}}\left[c_5\mathrm{sin}{\displaystyle \frac{\sqrt{6}\tau }{36M}}c_6\mathrm{cos}{\displaystyle \frac{\sqrt{6}\tau }{36M}}\right],`$
$`m`$ $`=`$ $`c_m,`$
$`S_{\widehat{t}\widehat{r}}`$ $`=`$ $`c_7\tau +c_8,`$
$`S_{\widehat{r}\widehat{\varphi }}`$ $`=`$ $`\pm {\displaystyle \frac{1}{2}}[c_7\tau +c_8]+c_1,`$
$`S_{\widehat{t}\widehat{\varphi }}`$ $`=`$ $`{\displaystyle \frac{\sqrt{2}\tau }{12M}}\left[{\displaystyle \frac{c_7}{2}}\tau +c_8\right]+c_2;`$ (5.2)
3. $`r>6M`$:
$`S_{\widehat{\theta }\widehat{\varphi }}`$ $`=`$ $`c_3\mathrm{cos}\omega _4\tau +c_4\mathrm{sin}\omega _4\tau ,`$
$`S_{\widehat{t}\widehat{\theta }}`$ $`=`$ $`c_5\mathrm{cos}\omega _3\tau +c_6\mathrm{sin}\omega _3\tau \pm {\displaystyle \frac{S_{\widehat{\theta }\widehat{\varphi }}}{\nu _K}},`$
$`S_{\widehat{r}\widehat{\theta }}`$ $`=`$ $`\gamma _K\nu _K\left[c_5\mathrm{sin}\omega _3\tau c_6\mathrm{cos}\omega _3\tau \right],`$
$`m`$ $`=`$ $`c_m,`$
$`S_{\widehat{t}\widehat{r}}`$ $`=`$ $`c_7\mathrm{cos}\omega _2\tau +c_8\mathrm{sin}\omega _2\tau \pm 2{\displaystyle \frac{\gamma _K}{\zeta _K}}{\displaystyle \frac{c_2}{43\gamma _K^2}},`$
$`S_{\widehat{r}\widehat{\varphi }}`$ $`=`$ $`\pm \nu _K[c_7\mathrm{cos}\omega _2\tau +c_8\mathrm{sin}\omega _2\tau ]+c_1+2{\displaystyle \frac{\nu _K}{\zeta _K}}\gamma _K{\displaystyle \frac{c_2}{43\gamma _K^2}},`$
$`S_{\widehat{t}\widehat{\varphi }}`$ $`=`$ $`\pm {\displaystyle \frac{2}{\gamma _K}}(43\gamma _K^2)^{1/2}[c_8\mathrm{cos}\omega _2\tau c_7\mathrm{sin}\omega _2\tau ]+c_93\gamma _K^2{\displaystyle \frac{c_2}{43\gamma _K^2}}\tau ,`$ (5.3)
where $`c_m,c_1,\mathrm{},c_9`$ are integration constants, and three real frequencies are defined for each open interval of radial values by
$`\omega _2=i\overline{\omega }_2=\zeta _K(43\gamma _K^2)^{1/2}=\sqrt{{\displaystyle \frac{M(r6M)}{r^3(r3M)}}},\omega _3=\zeta _K=\left({\displaystyle \frac{M}{r^3}}\right)^{1/2},`$
$`\omega _4=i\overline{\omega }_4=\zeta _K(3\gamma _K^22)^{1/2}={\displaystyle \frac{1}{r}}\sqrt{{\displaystyle \frac{M}{r3M}}}.`$ (5.4)
Consider first the open interval cases $`r6M`$. From Eq. (3), the total 4-momentum $`P`$ is given by
$`P`$ $`=`$ $`\left\{m\gamma _K\zeta _K\left[S_{\widehat{r}\widehat{\varphi }}(12\nu _K^2)\gamma _K^2c_12{\displaystyle \frac{\nu _K}{\gamma _K\zeta _K}}{\displaystyle \frac{c_2}{14\nu _K^2}}\right]\right\}e_{\widehat{t}}`$
$`{\displaystyle \frac{\gamma _K^2\zeta _K}{2}}\left\{(1+2\nu _K^2)\left[S_{\widehat{t}\widehat{\varphi }}+3{\displaystyle \frac{c_2}{14\nu _K^2}}\tau \right]+(14\nu _K^2)c_9\right\}e_{\widehat{r}}`$
$`+\left\{{\displaystyle \frac{\zeta _K}{\nu _K}}S_{\widehat{r}\widehat{\theta }}(1+2\nu _K^2)^{1/2}\left[c_3\mathrm{sin}\omega _4\tau c_4\mathrm{cos}\omega _4\tau \right]\right\}e_{\widehat{\theta }}`$
$`\pm \left\{m\gamma _K\nu _K{\displaystyle \frac{\zeta _K}{\nu _K}}\left[S_{\widehat{r}\widehat{\varphi }}(12\nu _K^2)\gamma _K^2c_12{\displaystyle \frac{\nu _K}{\gamma _K}}{\displaystyle \frac{c_2}{14\nu _K^2}}\right]\right\}e_{\widehat{\varphi }}.`$ (5.5)
We next impose the standard supplementary conditions. The CP conditions imply that $`S_{\widehat{t}\widehat{a}}=0`$, namely
$$c_2=c_3=c_4=c_5=c_6=c_7=c_8=c_9=0,$$
(5.6)
so that the only nonvanishing component of the spin tensor is $`S^{\widehat{z}}=S_{\widehat{r}\widehat{\varphi }}=c_1s`$, leaving arbitrary values for $`c_m`$ as well. From Eq. (5), the total 4-momentum $`P`$ becomes (using $`m_s=s\gamma \nu _K\zeta _K`$ which follows from Eq. (3.6))
$$P=mU_\pm +s\gamma \nu _K\zeta _KE_{\widehat{\varphi }},$$
(5.7)
with $`U_\pm `$ given by Eq. (2.7). Re-examing Eq. (3.7) shows that the spin-curvature force then acts radially, balancing the radial spin-orbit force.
The P conditions imply
$$S_{\widehat{t}\widehat{\varphi }}=0,S_{\widehat{r}\widehat{t}}\pm \nu _KS_{\widehat{r}\widehat{\varphi }}=0,S_{\widehat{\theta }\widehat{t}}\pm \nu _KS_{\widehat{\theta }\widehat{\varphi }}=0,$$
(5.8)
which lead only to the trivial solution
$$c_1=c_2=c_3=c_4=c_5=c_6=c_7=c_8=c_9=0,$$
(5.9)
with $`c_m`$ arbitrary, or in other words the components of the spin tensor must all be zero, which means that a non-zero spin is incompatible with geodesic motion for a spinning particle.
The T supplementary conditions when $`\nu =\pm \nu _K`$ imply
$`0`$ $`=`$ $`S_{\widehat{t}\widehat{\theta }}{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{\theta }}}{\mathrm{d}\tau _U}}+S_{\widehat{t}\widehat{r}}{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{r}}}{\mathrm{d}\tau _U}}+\gamma _K^2S_{\widehat{t}\widehat{\varphi }}{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{\varphi }}}{\mathrm{d}\tau _U}}\pm m\gamma _K^2\nu _KS_{\widehat{t}\widehat{\varphi }}`$
$`\gamma _K\zeta _K\{[S_{\widehat{t}\widehat{r}}S_{\widehat{t}\widehat{\varphi }}\pm \nu _KS_{\widehat{t}\widehat{\theta }}S_{\widehat{r}\widehat{\varphi }}]\gamma _K^2S_{\widehat{t}\widehat{\varphi }}[S_{\widehat{t}\widehat{r}}\nu _KS_{\widehat{r}\widehat{\varphi }}]\},`$
$`0`$ $`=`$ $`S_{\widehat{r}\widehat{\theta }}{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{\theta }}}{\mathrm{d}\tau _U}}\gamma _K^2[\pm \nu _KS_{\widehat{t}\widehat{r}}S_{\widehat{r}\widehat{\varphi }}]{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{\varphi }}}{\mathrm{d}\tau _U}}+\gamma _K{\displaystyle \frac{\zeta _K}{\nu _K}}[S_{\widehat{t}\widehat{r}}^2\nu _K^2S_{\widehat{r}\widehat{\theta }}^2]`$
$`\gamma _K^3{\displaystyle \frac{\zeta _K}{\nu _K}}\left[S_{\widehat{t}\widehat{r}}^2\nu _K(1+\nu _K^2)S_{\widehat{t}\widehat{r}}S_{\widehat{r}\widehat{\varphi }}+\nu _K^2S_{\widehat{r}\widehat{\varphi }}^2\right]m\gamma _K^2[S_{\widehat{t}\widehat{r}}\nu _KS_{\widehat{r}\widehat{\varphi }}],`$
$`0`$ $`=`$ $`\gamma _K^2[\pm \nu _KS_{\widehat{t}\widehat{\theta }}S_{\widehat{\theta }\widehat{\varphi }}]\left[{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{\varphi }}}{\mathrm{d}\tau _U}}\gamma _K\nu _K\zeta _KS_{\widehat{r}\widehat{\varphi }}\right]+S_{\widehat{r}\widehat{\theta }}{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{r}}}{\mathrm{d}\tau _U}}`$
$`+\gamma _K^2[S_{\widehat{t}\widehat{\theta }}\nu _KS_{\widehat{\theta }\widehat{\varphi }}]\left[m+\gamma _K{\displaystyle \frac{\zeta _K}{\nu _K}}S_{\widehat{t}\widehat{r}}\right]\gamma _K{\displaystyle \frac{\zeta _K}{\nu _K}}[S_{\widehat{t}\widehat{r}}S_{\widehat{t}\widehat{\theta }}\pm \nu _KS_{\widehat{r}\widehat{\theta }}S_{\widehat{t}\widehat{\varphi }}],`$
$`0`$ $`=`$ $`\pm \gamma _K^2\nu _KS_{\widehat{t}\widehat{\varphi }}{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{\varphi }}}{\mathrm{d}\tau _U}}+S_{\widehat{\theta }\widehat{\varphi }}{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{\theta }}}{\mathrm{d}\tau _U}}+S_{\widehat{r}\widehat{\varphi }}{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{r}}}{\mathrm{d}\tau _U}}+\gamma _K^3{\displaystyle \frac{\zeta _K}{\nu _K}}S_{\widehat{t}\widehat{\varphi }}[S_{\widehat{t}\widehat{r}}\nu _K^3S_{\widehat{r}\widehat{\varphi }}]`$
$`+m\gamma _K^2S_{\widehat{t}\widehat{\varphi }}\gamma _K{\displaystyle \frac{\zeta _K}{\nu _K}}[\nu _K^2S_{\widehat{r}\widehat{\theta }}S_{\widehat{\theta }\widehat{\varphi }}+S_{\widehat{t}\widehat{\varphi }}(S_{\widehat{t}\widehat{r}}\pm \nu _KS_{\widehat{r}\widehat{\varphi }})].`$ (5.10)
By substituting into these equations the solutions given by Eqs. (1)โ(3), we obtain the conditions
$$c_2=c_3=c_4=c_5=c_6=c_7=c_8=0,$$
(5.11)
implying that the only nonvanishing components of the spin tensor are
$$S^{\widehat{z}}=S_{\widehat{r}\widehat{\varphi }}=c_1,S_{\widehat{t}\widehat{\varphi }}=c_9,$$
(5.12)
and either
$$c_1,c_9\text{arbitrary},c_m=\pm \gamma _K\zeta _Kc_1,$$
(5.13)
which implies that spin component $`S^{\widehat{z}}`$ is proportional to the mass, locking them together by a constant of proportionality depending on the orbit velocity, or
$$c_1=0=c_9,c_m\text{arbitrary},$$
(5.14)
corresponding to the zero spin case where geodesic motion is of course allowed. In the former case the total spin invariant (2.14) reduces to
$$s^2=c_9^2+c_1^2,$$
(5.15)
so that the condition $`|s|/(mM)1`$ preserving the validity of the Mathisson-Papapetrou model reads
$$\frac{|s|}{mM}=\frac{1}{M\gamma _K\zeta _K}\left(1\frac{c_9^2}{c_1^2}\right)^{1/2}1,$$
(5.16)
implying either $`c_1c_9`$ or $`r3M`$ (where $`\gamma _K\mathrm{}`$). In the limit $`r3M`$ where the circular geodesics become null and require a separate treatment, one has a solution for which this spin component $`S^{\widehat{z}}`$ is fixed to have a value determined by the constant mass $`m`$ and the azimuthal velocity, the $`t`$-$`\varphi `$ component of the spin is arbitrary. If one takes the limit $`m0`$, then the component of the spin vector out of the orbit vanishes, leaving the spin vector locked to the direction of motion as found by who discussed the null geodesic case using the P supplementary conditions, the latter being the only physically relevant in such limit.
Finally consider the remaining case $`r=6M`$. Eq. (3) then shows that the total 4-momentum $`P`$ is given by
$`P`$ $`=`$ $`\left\{{\displaystyle \frac{2}{3}}\sqrt{3}m{\displaystyle \frac{\sqrt{6}}{108M}}\left[3S_{\widehat{r}\widehat{\varphi }}2c_1\right]\right\}e_{\widehat{t}}+\left\{{\displaystyle \frac{\sqrt{6}}{36M}}S_{\widehat{t}\widehat{\varphi }}+{\displaystyle \frac{\sqrt{3}}{2}}c_7\right\}e_{\widehat{r}}`$
$`\left\{{\displaystyle \frac{\sqrt{6}}{18M}}S_{\widehat{r}\widehat{\theta }}\pm {\displaystyle \frac{1}{6M}}\left[c_3\mathrm{sin}{\displaystyle \frac{\sqrt{3}\tau }{18M}}c_4\mathrm{cos}{\displaystyle \frac{\sqrt{3}\tau }{18M}}\right]\right\}e_{\widehat{\theta }}`$
$`\pm {\displaystyle \frac{1}{2}}\left\{{\displaystyle \frac{2}{3}}\sqrt{3}m{\displaystyle \frac{\sqrt{6}}{108M}}\left[3S_{\widehat{r}\widehat{\varphi }}2c_1\right]\right\}e_{\widehat{\varphi }}.`$ (5.17)
Imposing the standard supplementary conditions gives rise to the same result as for the general case $`r6M`$.
The CP conditions imply
$$c_2=c_3=c_4=c_5=c_6=c_7=c_8=0,$$
(5.18)
so that the only nonvanishing component of the spin tensor is $`S^{\widehat{z}}=S_{\widehat{r}\widehat{\varphi }}=c_1`$, for arbitrary values of $`c_m`$, leading to constant mass $`m`$. The P conditions give only the trivial solution
$$c_1=c_2=c_3=c_4=c_5=c_6=c_7=c_8=0,$$
(5.19)
with $`c_m`$ arbitrary, leading to constant mass $`m`$. Finally, the T conditions imply
$$c_3=c_4=c_5=c_6=c_7=c_8=0,$$
(5.20)
so that the only nonvanishing components of the spin tensor are
$$S^{\widehat{z}}=S_{\widehat{r}\widehat{\varphi }}=c_1,S_{\widehat{t}\widehat{\varphi }}=c_2,$$
(5.21)
and either
$$c_1,c_2\text{arbitrary},c_m=\pm \frac{\sqrt{2}}{18M}c_1$$
(5.22)
or
$$c_1=0=c_2,c_m\text{arbitrary},$$
(5.23)
with constant mass $`m`$ in both cases. In the former case the spin invariant (2.14) reduces to
$$s^2=c_2^2+c_1^2,$$
(5.24)
so that the condition $`|s|/(mM)1`$ reads
$$\frac{|s|}{mM}=\frac{18}{\sqrt{2}}\left(1\frac{c_2^2}{c_1^2}\right)^{1/2}1,$$
(5.25)
implying $`c_1c_2`$.
Thus if the center of mass of the test particle is constrained to be a circular geodesic, either the spin is forced to be zero or have an arbitrary constant value of the single nonzero component $`S^{\widehat{z}}`$ of the spin vector out of the plane of the orbit.
## 6 The general case: $`\nu 0`$ and $`\nu \pm \nu _K`$
For general circular orbits excluding the previous cases $`\nu =0`$ and $`\nu =\pm \nu _K`$, the solutions of the equations of motion for the components of the spin tensor and the mass $`m`$ of the spinning particle are
$`S_{\widehat{\theta }\widehat{\varphi }}`$ $`=`$ $`A\mathrm{cos}\mathrm{\Omega }\tau +B\mathrm{sin}\mathrm{\Omega }\tau ,`$ (6.1)
$`S_{\widehat{t}\widehat{\theta }}`$ $`=`$ $`C\mathrm{cos}\mathrm{\Omega }_1\tau +D\mathrm{sin}\mathrm{\Omega }_1\tau +FS_{\widehat{\theta }\widehat{\varphi }},F={\displaystyle \frac{3\nu \nu _K^2}{[\nu ^2(1+2\nu _K^2)\nu _K^2(1\nu _K^2)]}},`$ (6.2)
$`S_{\widehat{r}\widehat{\theta }}`$ $`=`$ $`{\displaystyle \frac{\gamma \nu \zeta _K(1+2\nu _K^2)(\nu ^2\nu _K^2)}{\mathrm{\Omega }\nu _K[\nu ^2(1+2\nu _K^2)\nu _K^2(1\nu _K^2)]}}[A\mathrm{sin}\mathrm{\Omega }\tau B\mathrm{cos}\mathrm{\Omega }\tau ]`$
$`\nu _K\gamma _K[D\mathrm{cos}\mathrm{\Omega }_1\tau C\mathrm{sin}\mathrm{\Omega }_1\tau ],`$ (6.3)
$`S_{\widehat{r}\widehat{\varphi }}`$ $`=`$ $`{\displaystyle \frac{\nu _K}{\zeta _K}}{\displaystyle \frac{\nu ^2\nu _K^2}{\frac{\nu ^2}{\gamma _K^2}+\nu _K^2(2+\nu _K^2)}}\left[\gamma \nu c_m\gamma _K^2{\displaystyle \frac{\nu _K}{\zeta _K}}{\displaystyle \frac{\nu ^2(14\nu _K^2)+\nu _K^2(2+\nu _K^2)}{(\nu ^2\nu _K^2)^2}}c_0\right]`$
$`+c_1e^{\mathrm{\Omega }_+\tau }+c_2e^{\mathrm{\Omega }_+\tau }+c_3e^{\mathrm{\Omega }_{}\tau }+c_4e^{\mathrm{\Omega }_{}\tau },`$ (6.4)
$`S_{\widehat{t}\widehat{r}}`$ $`=`$ $`{\displaystyle \frac{\nu _K}{\zeta _K}}{\displaystyle \frac{1}{\frac{\nu ^2}{\gamma _K^2}+\nu _K^2(2+\nu _K^2)}}[\gamma (\nu ^2\nu _K^2)c_m+\nu {\displaystyle \frac{\nu _K}{\zeta _K}}c_0]+{\displaystyle \frac{1}{2\nu \nu _K^2(1+2\nu _K^2)}}`$
$`\{(3\nu _K^2+\mathrm{\Phi })\left[c_1e^{\mathrm{\Omega }_+\tau }+c_2e^{\mathrm{\Omega }_+\tau }\right]+(3\nu _K^2\mathrm{\Phi })\left[c_3e^{\mathrm{\Omega }_{}\tau }+c_4e^{\mathrm{\Omega }_{}\tau }\right]\},`$ (6.5)
$`S_{\widehat{t}\widehat{\varphi }}`$ $`=`$ $`{\displaystyle \frac{1}{\gamma }}{\displaystyle \frac{\nu _K}{\zeta _K}}{\displaystyle \frac{1}{\nu ^2\nu _K^2}}\{\mathrm{\Omega }_+[{\displaystyle \frac{3\nu _K^2+\mathrm{\Phi }}{2\nu _K^2(1+2\nu _K^2)}}1][c_1e^{\mathrm{\Omega }_+\tau }c_2e^{\mathrm{\Omega }_+\tau }]`$
$`+\mathrm{\Omega }_{}[{\displaystyle \frac{3\nu _K^2\mathrm{\Phi }}{2\nu _K^2(1+2\nu _K^2)}}1][c_3e^{\mathrm{\Omega }_{}\tau }c_4e^{\mathrm{\Omega }_{}\tau }]\},`$ (6.6)
$`m`$ $`=`$ $`\gamma {\displaystyle \frac{\zeta _K}{\nu _K}}[\nu S_{\widehat{r}\widehat{\varphi }}\nu _K^2S_{\widehat{t}\widehat{r}}]+c_m;`$ (6.7)
where $`A,B,C,D,c_m,c_0,\mathrm{},c_4`$ are integration constants, and the real positive frequencies $`\mathrm{\Omega }`$ and $`\mathrm{\Omega }_1`$ are given by
$$\mathrm{\Omega }=\gamma \zeta _K(1+2\nu _K^2)^{1/2}\frac{|\nu |}{\nu _K},\mathrm{\Omega }_1=\frac{\gamma \zeta _K}{\gamma _K},$$
(6.8)
assumed to be distinct for the above equations to be valid, and the remaining abbreviations are
$$\mathrm{\Omega }_\pm =\frac{\gamma }{\gamma _K}\frac{\zeta _K}{\nu _K}\left[\overline{\nu }^2\nu ^2\pm \frac{\gamma _K^2}{2}\mathrm{\Phi }\right]^{1/2},\mathrm{\Phi }=3\nu _K^2\left[1\frac{\nu ^2}{\stackrel{~}{\nu }^2}\right]^{1/2},$$
(6.9)
with
$$\overline{\nu }^2=\frac{\gamma _K^2\nu _K^2}{2}(1+2\nu _K^2),\stackrel{~}{\nu }^2=\frac{9}{8}\frac{\gamma _K^2\nu _K^2}{(1+2\nu _K^2)}.$$
(6.10)
The behaviors of the azimuthal velocities $`\overline{\nu }`$, $`\stackrel{~}{\nu }`$ and $`\nu _K`$ as functions of the radial parameter $`r/M`$ are compared in Fig. 1. They all coincide at $`r=6M`$, where $`\overline{\nu }=\stackrel{~}{\nu }=\nu _K=1/2`$; for $`2M<r<6M`$ it is $`\stackrel{~}{\nu }<\overline{\nu }`$, while $`\stackrel{~}{\nu }>\overline{\nu }`$ for $`r>6M`$.
The quantities $`\mathrm{\Omega }_\pm `$ also lead to angular velocities for certain intervals of values of the azimuthal velocity $`\nu `$. In fact we are interested in those values for which $`\mathrm{\Omega }_+`$ and/or $`\mathrm{\Omega }_{}`$ are purely imaginary, since the imaginary parts can be interpreted as additional frequencies characterizing spin precession. One must distinguish the cases $`2M<r<6M`$ and $`r>6M`$, referring to Fig. 1 and to Eq. (6.9):
* $`r>6M`$:
+ if $`\nu >\stackrel{~}{\nu }`$ (region I), the quantities $`\mathrm{\Omega }_\pm `$ are both complex;
+ if $`\nu =\stackrel{~}{\nu }`$, $`\mathrm{\Omega }_+=\mathrm{\Omega }_{}`$ is purely imaginary, since $`\stackrel{~}{\nu }>\overline{\nu }`$;
+ if $`\overline{\nu }<\nu <\stackrel{~}{\nu }`$ (region II), $`\mathrm{\Omega }_{}`$ is purely imaginary, while $`\mathrm{\Omega }_+`$ can be either real (even zero) or purely imaginary;
+ if $`\nu <\overline{\nu }`$ (region III), $`\mathrm{\Omega }_+`$ is purely imaginary, while $`\mathrm{\Omega }_{}`$ can be either real (even zero) or purely imaginary;
* $`2M<r<6M`$:
+ if $`\nu >\stackrel{~}{\nu }`$ (region IV), the quantities $`\mathrm{\Omega }_\pm `$ are both complex;
+ if $`\nu =\stackrel{~}{\nu }`$, $`\mathrm{\Omega }_+=\mathrm{\Omega }_{}`$ is real, since $`\stackrel{~}{\nu }<\overline{\nu }`$;
+ if $`\nu <\stackrel{~}{\nu }`$ (region V), $`\mathrm{\Omega }_+`$ is real, since $`\stackrel{~}{\nu }<\overline{\nu }`$, while $`\mathrm{\Omega }_{}`$ can be either real (even zero) or purely imaginary.
All of these remarks so far assume that the two frequencies $`\mathrm{\Omega }`$ and $`\mathrm{\Omega }_1`$ are distinct, necessary for the decoupling procedure which leads to this solution. A different result follows in the special case $`\mathrm{\Omega }=\mathrm{\Omega }_1`$. This occurs for the particular value of the azimuthal velocity
$$\nu _0=\pm \frac{\nu _K}{\gamma _K}(1+2\nu _K^2)^{1/2}=\pm \sqrt{\frac{M(r3M)}{r(r2M)}},$$
(6.11)
which vanishes at $`r=3M`$ and is real for $`r>3M`$, rising to its peak speed at $`r3.934M`$ and decreasing asymptotically towards the geodesic speed from below as $`r\mathrm{}`$. The solutions for the components $`S_{\widehat{\theta }\widehat{\varphi }}`$, $`S_{\widehat{t}\widehat{\theta }}`$ and $`S_{\widehat{r}\widehat{\theta }}`$ of the spin tensor are given by
$`S_{\widehat{\theta }\widehat{\varphi }}`$ $`=`$ $`A\mathrm{cos}\mathrm{\Omega }\tau +B\mathrm{sin}\mathrm{\Omega }\tau ,`$ (6.12)
$`S_{\widehat{t}\widehat{\theta }}`$ $`=`$ $`\left[C{\displaystyle \frac{3}{2}}{\displaystyle \frac{\gamma _0^2\nu _0}{\mathrm{\Omega }}}\zeta _K^2(AB\mathrm{\Omega }\tau )\right]\mathrm{cos}\mathrm{\Omega }\tau +\left[D{\displaystyle \frac{3}{2}}{\displaystyle \frac{\gamma _0^2\nu _0}{\mathrm{\Omega }^2}}\zeta _K^2A\tau \right]\mathrm{sin}\mathrm{\Omega }\tau ,`$ (6.13)
$`S_{\widehat{r}\widehat{\theta }}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{\zeta _K}{\nu _K}}{\displaystyle \frac{\gamma _0^3}{\mathrm{\Omega }^2}}\{[3\nu _0\nu _K^2\zeta _K^2(B+A\mathrm{\Omega }\tau )+2{\displaystyle \frac{\mathrm{\Omega }^2}{\gamma _0^2}}(\nu _0B\nu _K^2D)]\mathrm{cos}\mathrm{\Omega }\tau `$
$`+[3\nu _0\nu _K^2\zeta _K^2(2A+B\mathrm{\Omega }\tau )2{\displaystyle \frac{\mathrm{\Omega }^2}{\gamma _0^2}}(\nu _0A\nu _K^2C)]\mathrm{sin}\mathrm{\Omega }\tau \},`$ (6.14)
with
$$\mathrm{\Omega }\mathrm{\Omega }_1=\frac{\zeta _K}{\nu _K}\left[\frac{1+2\nu _K^2}{1+\nu _K^2(1+\nu _K^2)}\right]^{1/2}=\frac{\sqrt{M}}{r}\left[\frac{r3M}{r(r3M)+3M^2}\right]^{1/2},$$
(6.15)
while that corresponding to the remaining components as well as to the varying mass $`m`$ are obtained simply by evaluating the general solutions (6)โ(6.7) at $`\nu =\nu _0`$. The reality properties of the quantities $`\mathrm{\Omega }_\pm `$ can be determined as done in the general case, noting that $`\nu _0<\stackrel{~}{\nu }`$ (corresponding to region V) always holds in the interval $`3M<r<6M`$. For $`r>6M`$ however, we must distinguish two different regions, a) $`6M<r<\overline{r}_0`$, with $`\overline{r}_0=6M(1+\sqrt{2}/2)10.24M`$ such that $`\nu _0=\overline{\nu }`$, where $`\nu _0<\overline{\nu }`$ (corresponding to region III), and b) $`r>\overline{r}_0`$, where $`\nu _0>\overline{\nu }`$ (corresponding to region II).
The behavior of $`S`$, $`U`$ and $`P`$ along the world line itself is completely determined by the initial conditions
$`S_{\widehat{\alpha }\widehat{\beta }}(0),{\displaystyle \frac{\mathrm{d}S_{\widehat{\alpha }\widehat{\beta }}}{\mathrm{d}\tau _U}}|_{\tau =0},`$ (6.16)
and the corresponding conditions on the mass $`m`$ of the particle which follow from Eq. (6.7). Thus in the special case in which the โcenter of mass lineโ is directed along a circular orbit, the completion of the scheme for the spinning test particle is equivalent to a choice of initial conditions.
In principle the components of the spin tensor which are not constants should precess with the different frequencies which appear in Eqs. (6.1)โ(6), leading to non-periodic motion, a feature that seems to characterize the general situation in the Schwarzschild and Kerr spacetimes. However, this does not occur in practice once the CP, P and T supplementary conditions are imposed, as we will see below. It turns out that the nonvanishing components of the spin tensor are all constant in the CP and T cases, while the motion is periodic with a unique frequency in the P case. As one might expect, the particle mass $`m`$ turns out to be constant in all three cases.
### 6.1 The CP supplementary conditions
The CP supplementary conditions require
$$S_{\widehat{t}\widehat{r}}=0,S_{\widehat{t}\widehat{\theta }}=0,S_{\widehat{t}\widehat{\varphi }}=0.$$
(6.17)
From Eq. (6) this forces
$$c_1=c_2=c_3=c_4=0,c_0=\frac{\gamma }{\nu }\frac{\zeta _K}{\nu _K}(\nu ^2\nu _K^2)c_m.$$
(6.18)
Substituting these values into Eq. (6) then gives
$$c_m=\frac{\gamma \nu }{\gamma _K^2}\frac{\zeta _K}{\nu _K}S_{\widehat{r}\widehat{\varphi }},$$
(6.19)
so from Eq. (6.7) we get
$$S_{\widehat{r}\widehat{\varphi }}=s=\frac{m}{\gamma \nu }\frac{1}{\nu _K\zeta _K}.$$
(6.20)
Finally, from Eqs. (6.1) and (6) it follows that
$$S_{\widehat{\theta }\widehat{\varphi }}=0,S_{\widehat{r}\widehat{\theta }}=0.$$
(6.21)
However, from Eq. (3) it follows that
$$P=s\nu _K\zeta _Ke_{\widehat{\varphi }},$$
(6.22)
since $`m_s=s\gamma \nu _K\zeta _K=m/\nu `$, a consequence of Eqs. (3.6) and (6.20). This result is unphysical since the total 4-momentum $`P`$ is spacelike.
### 6.2 The P supplementary conditions
The P supplementary conditions require
$$S_{\widehat{t}\widehat{\varphi }}=0,S_{\widehat{r}\widehat{t}}+S_{\widehat{r}\widehat{\varphi }}\nu =0,S_{\widehat{\theta }\widehat{t}}+S_{\widehat{\theta }\widehat{\varphi }}\nu =0.$$
(6.23)
Under these conditions the components of the spin vector $`S_U`$ in the local rest space of the particle, $`S_U^\beta =\frac{1}{2}\eta _\alpha {}_{}{}^{\beta \gamma \delta }U_{}^{\alpha }S_{\gamma \delta }`$, expressed in the Frenet-Serret frame, are just
$$(S_U^1,S_U^2,S_U^3)=(\gamma ^1S_{\widehat{\theta }\widehat{\varphi }},S_{\widehat{r}\widehat{\theta }},\gamma ^1S_{\widehat{r}\widehat{\varphi }}).$$
(6.24)
Comparing the first Eq. (6.23) with Eq. (6) we get
$$c_1=c_2=c_3=c_4=0,$$
(6.25)
so that $`S_{\widehat{t}\widehat{r}}`$, $`S_{\widehat{r}\widehat{\varphi }}`$ and the particle mass $`m`$ are all constant. Eqs. (6) and (6) together with the second of the Pirani conditions Eq. (6.23) imply
$$c_0=\frac{\gamma }{\nu }\frac{\zeta _K}{\nu _K}\frac{1+\nu ^2}{\gamma _K^2}\frac{(\nu ^2\nu _K^2)^2}{\nu ^2(25\nu _K^2)+\nu _K^2(1+2\nu _K^2)}c_m,$$
(6.26)
hence from Eq. (6.7)
$$c_m=\left[1+\frac{1}{\gamma _K^2}\frac{\nu ^2\nu _K^2}{\nu ^2(14\nu _K^2)+\nu _K^2(2+\nu _K^2)}\right]m.$$
(6.27)
Next by substituting these values of the constants $`c_0`$ and $`c_m`$ into Eq. (6), we obtain
$$\gamma S_U^3=S^{\widehat{z}}=S_{\widehat{r}\widehat{\varphi }}=m\frac{\gamma }{\nu }\frac{\nu _K}{\zeta _K}\frac{\nu ^2\nu _K^2}{\frac{\gamma ^2}{\gamma _K^2}(\nu ^2\nu _K^2)+3\nu _K^2}.$$
(6.28)
Finally comparing the last of the Pirani conditions Eq. (6.23) with Eqs. (6.1)โ(6.2) and Eqs. (6.12)โ(6.13) leads to two possibilities: either
$$\text{case P1: }A=B=C=D=0,$$
(6.29)
which places no constraint on $`\nu `$ and the spin vector is constant and out of the plane of the orbit, or
$$\text{case P2: }C=0=D,F=\nu ,$$
(6.30)
the latter of which (again from Eq. (6.2)) leads to the special azimuthal velocity
$$\nu =\nu _{(P2)}=\pm 2\nu _K\left(\frac{1\nu _K^2/4}{1+2\nu _K^2}\right)^{1/2}=\pm 2\left(\frac{M(r9M/4)}{r(r2M)}\right)^{1/2}.$$
(6.31)
The case P1 has been already considered previously , leaving only P2 to be considered here. The corresponding azimuthal speed $`\nu _{(P2)}`$ vanishes at $`r=9/4M`$ and is real for $`r>9/4M`$, rising to a maximum speed of 1 at $`r=3M`$, corresponding to the two null circular geodesics, and decreasing asymptotically towards twice the geodesic speed from below as $`r\mathrm{}`$.
The corresponding values of $`\gamma `$ and $`\mathrm{\Omega }`$ are respectively
$$\gamma _{(P2)}=\frac{(1+2\nu _K^2)^{1/2}}{1\nu _K^2}=\frac{\sqrt{r(r2M)}}{r3M},$$
(6.32)
and
$$\mathrm{\Omega }_{(P2)}=\gamma _K^2\zeta _K[(4\nu _K^2)(1+2\nu _K^2)]^{1/2}=\frac{\sqrt{M}}{r}\frac{\sqrt{4r9M}}{r3M},$$
(6.33)
using Eq. (2.6). To get the anglular velocity of precession with respect to a frame which is nonrotating with respect to infinity, one must subtract away the precession angular velocity $`\mathrm{\Omega }_U=\gamma \nu /r`$ of the spherical frame. In the case P2 one finds $`\mathrm{\Omega }_{(P2)}\mathrm{\Omega }_U=0`$, so the spin does not precess with respect to a frame which is nonrotating at infinity.
Substituting these values back into Eq. (6.28) then leads to
$$S_U^3=\pm \frac{m}{2\mathrm{\Omega }_{(P2)}}.$$
(6.34)
The remaining nonzero spin components (6.1)โ(6) can then be expressed in the form
$`S_U^1=\gamma ^1S_{\widehat{\theta }\widehat{\varphi }}`$ $`=`$ $`\gamma ^1[A\mathrm{cos}\mathrm{\Omega }_{(P2)}\tau _U+B\mathrm{sin}\mathrm{\Omega }_{(P2)}\tau _U],`$
$`S_U^2=S_{\widehat{r}\widehat{\theta }}`$ $`=`$ $`\gamma ^1[A\mathrm{sin}\mathrm{\Omega }_{(P2)}\tau _UB\mathrm{cos}\mathrm{\Omega }_{(P2)}\tau _U],`$ (6.35)
$`S_U^1(0)=\gamma ^1A,S_U^2(0)=\gamma ^1B,S_U^3(0)=\pm {\displaystyle \frac{m}{2\mathrm{\Omega }_{(P2)}}},`$ (6.36)
leading to
$$\left(\begin{array}{c}S_U^1\\ S_U^2\\ S_U^3\end{array}\right)=\left(\begin{array}{ccc}\mathrm{cos}\varphi & \mathrm{sin}\varphi & 0\\ \mathrm{sin}\varphi & \mathrm{cos}\varphi & 0\\ 0& 0& 1\end{array}\right)\left(\begin{array}{c}S_U^1(0)\\ S_U^2(0)\\ S_U^3(0)\end{array}\right).$$
(6.37)
The spin invariant (2.14) becomes in this case
$$s^2=\frac{1}{\gamma ^2}\left[A^2+B^2+\frac{m^2}{4\zeta _K^2}\frac{1}{4\nu _K^2}\right].$$
(6.38)
The Mathisson-Papapetrou model is valid if the condition $`|s|/(mM)1`$ is satisfied. From the previous equation we have that either $`\gamma \mathrm{}`$ or the sum of the bracketed terms must be small, i.e.,
$$\frac{A^2}{m^2M^2}1,\frac{B^2}{m^2M^2}1,[4M^2\zeta _K^2(4\nu _K^2)]^11.$$
(6.39)
The latter possibility cannot occur for any allowed values of $`r/M`$, since the third term (which is dimensionless) of (6.39) is always greater than $`1.88`$, as is easily verified. The former possibility is realized only in the case of ultrarelativistic motion, which Eq. (6.32) implies occurs only as $`r3M`$, where the orbits approach null geodesics and the limit $`m0`$ forces the component of the spin vector out of the plane of the orbit to vanish, locking the spin vector to the direction of motion exactly as discussed by Mashhoon .
It is well known that the spin vector $`S_U=S_U^iE_i`$ lying in the local rest space of $`U`$ is Fermi-Walker transported along $`U`$ in the P case, so it must satisfy
$$0=\frac{D_{(\mathrm{fw})}S_U}{\mathrm{d}\tau _U}P(U)\frac{DS_U}{\mathrm{d}\tau _U}=\left[\frac{\mathrm{d}S_U^1}{\mathrm{d}\tau _U}+S_U^2\tau _1\right]E_1+\left[\frac{\mathrm{d}S_U^2}{\mathrm{d}\tau _U}S_U^1\tau _1\right]E_2,$$
(6.40)
from (2.10), where $`P(U)_\alpha ^\mu =\delta _\alpha ^\mu +U^\mu U_\alpha `$ projects into the local rest space of $`U`$. To check this we must show that the following two equations are identically satisfied
$$\frac{\mathrm{d}S_U^1}{\mathrm{d}\tau _U}+\tau _1S_U^2=0,\frac{\mathrm{d}S_U^2}{\mathrm{d}\tau _U}\tau _1S_U^1=0.$$
(6.41)
But these two equations follow immediately from (6.2), since $`\tau _1=\mathrm{\Omega }_{(P2)}`$ results from the direct evaluation of the expression (2) for $`\tau _1`$, with $`\nu `$ given by (6.31).
Thus given the rest mass $`m`$ of the test particle, the constant component of the spin orthogonal to the orbit is then fixed by the orbit parameters, while the component in the plane of the orbit as seen within the local rest space of the particle itself is locked to a direction which is fixed with respect to the distant observers, since the angle of precession with respect to the spherical axes is exactly the azimuthal angle of the orbit, but in the opposite sense. In other words the precession of the spin, which introduces a time varying force into the mix, must be locked to the first torsion of the orbit itself in order to maintain the alignment of the 4-velocity with a static direction in the spacetime, and the spin does not precess with respect to observers at spatial infinity. Furthermore, the specific spin of the test particle cannot be made arbitrarily small except near the limiting radius where the 4-velocity of this solution goes null, and the spin vector is then locked to the direction of motion. Apparently the imposition of a circular orbit on the center of mass world line of the test particle is just too strong a condition to describe an interesting spin precession.
The total 4-momentum $`P`$ given by Eqs. (3) and (3.6) can be written in this case as
$`P`$ $`=`$ $`mU+m_sE_2+\gamma ^1\left[{\displaystyle \frac{\mathrm{d}S_U^2}{\mathrm{d}\tau _U}}\gamma \nu _K\zeta _KS_{\widehat{r}\widehat{\theta }}\right]e_{\widehat{\theta }}`$ (6.42)
$`=`$ $`mU+m_sE_2\left({\displaystyle \frac{\gamma \nu ^2}{r}}+\nu _K\zeta _K\right)S_U^2e_{\widehat{z}},`$
with $`\nu `$ given by (6.31) and $`m_s`$ a constant
$$m_s=\gamma \frac{\zeta _K}{\nu _K}(\nu S_{\widehat{t}\widehat{r}}\nu _K^2S_{\widehat{r}\widehat{\varphi }})=\gamma ^2\frac{\zeta _K}{\nu _K}(\nu ^2\nu _K^2)S_U^3,$$
(6.43)
but the final term in $`P`$ (out of the plane of the orbit) oscillates as the spin precesses in the plane of the orbit. Note that the radial component of $`P`$ is zero.
The spin-curvature force (3.7) simplifies to
$`F^{(\mathrm{sc})}`$ $`=`$ $`\gamma \zeta _K^2\left\{\left[2S_{\widehat{t}\widehat{r}}+\nu S_{\widehat{r}\widehat{\varphi }}\right]e_{\widehat{r}}\left[S_{\widehat{t}\widehat{\theta }}+2\nu S_{\widehat{\theta }\widehat{\varphi }}\right]e_{\widehat{\theta }}\right\}`$ (6.44)
$`=`$ $`3\gamma ^2\nu \zeta _K^2(S_U^2e_{\widehat{r}}S_U^1e_{\widehat{\theta }}),`$
while the term on the left hand side of Eq. (1.1) can be written as
$`{\displaystyle \frac{DP}{\mathrm{d}\tau _U}}`$ $`=`$ $`[m\kappa m_s\tau _1]e_{\widehat{r}}\gamma \zeta _K^2\left[S_{\widehat{t}\widehat{\theta }}+2\nu S_{\widehat{\theta }\widehat{\varphi }}\right]e_{\widehat{\theta }}`$ (6.45)
$`=`$ $`[m\kappa m_s\tau _1]e_{\widehat{r}}3\gamma ^2\nu \zeta _K^2S_U^1e_{\widehat{\theta }}.`$
The force balance equation (3.8) reduces to
$`ma(U)_{\widehat{r}}`$ $`=`$ $`F_{\widehat{r}}^{(\mathrm{so})}+F_{\widehat{r}}^{(\mathrm{sc})},`$
$`ma(U)_{\widehat{\theta }}`$ $`=`$ $`0=F_{\widehat{\theta }}^{(\mathrm{so})}+F_{\widehat{\theta }}^{(\mathrm{sc})},`$ (6.46)
where
$`ma(U)_{\widehat{r}}=m\kappa ,F_{\widehat{r}}^{(\mathrm{so})}=m_s\left({\displaystyle \frac{DE_{\widehat{\varphi }}}{d\tau _U}}\right)_{\widehat{r}}=m_s\tau _1,`$ (6.47)
$`F_{\widehat{r}}^{(\mathrm{sc})}=3\gamma ^2\nu \zeta _K^2S_U^3,F_{\widehat{\theta }}^{(\mathrm{sc})}=F_{\widehat{\theta }}^{(\mathrm{so})}=3\gamma ^2\nu \zeta _K^2S_U^1.`$
### 6.3 The Tulczyjew (T) supplementary conditions
The T supplementary conditions imply from (3.1)
$`0`$ $`=`$ $`S_{\widehat{t}\widehat{\theta }}{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{\theta }}}{\mathrm{d}\tau _U}}+S_{\widehat{t}\widehat{r}}{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{r}}}{\mathrm{d}\tau _U}}+\gamma ^2S_{\widehat{t}\widehat{\varphi }}{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{\varphi }}}{\mathrm{d}\tau _U}}+m\gamma ^2\nu S_{\widehat{t}\widehat{\varphi }}`$
$`\gamma {\displaystyle \frac{\zeta _K}{\nu _K}}\{[\nu S_{\widehat{t}\widehat{r}}S_{\widehat{t}\widehat{\varphi }}+\nu _K^2S_{\widehat{t}\widehat{\theta }}S_{\widehat{r}\widehat{\varphi }}]\gamma ^2S_{\widehat{t}\widehat{\varphi }}[\nu S_{\widehat{t}\widehat{r}}\nu _K^2S_{\widehat{r}\widehat{\varphi }}]\},`$
$`0`$ $`=`$ $`S_{\widehat{r}\widehat{\theta }}{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{\theta }}}{\mathrm{d}\tau _U}}\gamma ^2[\nu S_{\widehat{t}\widehat{r}}S_{\widehat{r}\widehat{\varphi }}]{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{\varphi }}}{\mathrm{d}\tau _U}}+\gamma {\displaystyle \frac{\zeta _K}{\nu _K}}[S_{\widehat{t}\widehat{r}}^2\nu _K^2S_{\widehat{r}\widehat{\theta }}^2]`$
$`\gamma ^3{\displaystyle \frac{\zeta _K}{\nu _K}}\left[S_{\widehat{t}\widehat{r}}^2\nu (1+\nu _K^2)S_{\widehat{t}\widehat{r}}S_{\widehat{r}\widehat{\varphi }}+\nu _K^2S_{\widehat{r}\widehat{\varphi }}^2\right]m\gamma ^2[S_{\widehat{t}\widehat{r}}\nu S_{\widehat{r}\widehat{\varphi }}],`$
$`0`$ $`=`$ $`\gamma ^2[\nu S_{\widehat{t}\widehat{\theta }}S_{\widehat{\theta }\widehat{\varphi }}]\left[{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{\varphi }}}{\mathrm{d}\tau _U}}\gamma \nu _K\zeta _KS_{\widehat{r}\widehat{\varphi }}\right]+S_{\widehat{r}\widehat{\theta }}{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{r}}}{\mathrm{d}\tau _U}}`$
$`+\gamma ^2[S_{\widehat{t}\widehat{\theta }}\nu S_{\widehat{\theta }\widehat{\varphi }}]\left[m+\gamma {\displaystyle \frac{\zeta _K}{\nu _K}}S_{\widehat{t}\widehat{r}}\right]\gamma {\displaystyle \frac{\zeta _K}{\nu _K}}[S_{\widehat{t}\widehat{r}}S_{\widehat{t}\widehat{\theta }}+\nu S_{\widehat{r}\widehat{\theta }}S_{\widehat{t}\widehat{\varphi }}],`$
$`0`$ $`=`$ $`\gamma ^2\nu S_{\widehat{t}\widehat{\varphi }}{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{\varphi }}}{\mathrm{d}\tau _U}}+S_{\widehat{\theta }\widehat{\varphi }}{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{\theta }}}{\mathrm{d}\tau _U}}+S_{\widehat{r}\widehat{\varphi }}{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{r}}}{\mathrm{d}\tau _U}}+\gamma ^3{\displaystyle \frac{\zeta _K}{\nu _K}}S_{\widehat{t}\widehat{\varphi }}[S_{\widehat{t}\widehat{r}}\nu \nu _K^2S_{\widehat{r}\widehat{\varphi }}]`$
$`+m\gamma ^2S_{\widehat{t}\widehat{\varphi }}\gamma {\displaystyle \frac{\zeta _K}{\nu _K}}[\nu _K^2S_{\widehat{r}\widehat{\theta }}S_{\widehat{\theta }\widehat{\varphi }}+S_{\widehat{t}\widehat{\varphi }}(S_{\widehat{t}\widehat{r}}+\nu S_{\widehat{r}\widehat{\varphi }})].`$ (6.48)
By solving for the first derivatives, a straightforward calculation shows that the above set of equations simplifies to
$`0`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{r}}}{\mathrm{d}\tau _U}}+m{\displaystyle \frac{S_{\widehat{t}\widehat{\theta }}}{S_{\widehat{r}\widehat{\theta }}}}\gamma \nu {\displaystyle \frac{\zeta _K}{\nu _K}}S_{\widehat{t}\widehat{\varphi }},`$ (6.49)
$`0`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{\varphi }}}{\mathrm{d}\tau _U}}+m\nu +\gamma {\displaystyle \frac{\zeta _K}{\nu _K}}[\nu S_{\widehat{t}\widehat{r}}\nu _K^2S_{\widehat{r}\widehat{\varphi }}],`$ (6.50)
$`0`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{\theta }}}{\mathrm{d}\tau _U}}m{\displaystyle \frac{S_{\widehat{t}\widehat{r}}}{S_{\widehat{r}\widehat{\theta }}}}\gamma \nu _K\zeta _KS_{\widehat{r}\widehat{\theta }},`$ (6.51)
$`0`$ $`=`$ $`{\displaystyle \frac{m}{\gamma S_{\widehat{r}\widehat{\theta }}}}[S_{\widehat{t}\widehat{\theta }}S_{\widehat{r}\widehat{\varphi }}S_{\widehat{r}\widehat{\theta }}S_{\widehat{t}\widehat{\varphi }}S_{\widehat{t}\widehat{r}}S_{\widehat{\theta }\widehat{\varphi }}],`$ (6.52)
provided that $`S_{\widehat{r}\widehat{\theta }}0`$ is assumed. Substituting Eqs. (1.30), (6.7) and then Eqs. (A.3), (A.3) into Eq. (6.50) (see the equations listed in A.3) leads to
$$0=S_{\widehat{r}\widehat{\varphi }}\frac{\nu _K}{\zeta _K}\left[\gamma \nu c_m\frac{\nu _K}{\zeta _K}\frac{c_0}{\nu ^2\nu _K^2}\right].$$
(6.53)
Substituting Eq. (6) into Eq. (6.53) leads to
$$c_1=c_2=c_3=c_4=0,c_0=\frac{\gamma \nu }{\gamma _K^2}\frac{\zeta _K}{\nu _K}\left[1+\frac{1}{\gamma ^2}\frac{1}{2+\nu _K^2}\right]c_m,$$
(6.54)
implying that $`S_{\widehat{t}\widehat{\varphi }}=0`$, and $`S_{\widehat{t}\widehat{r}}`$, $`S_{\widehat{r}\widehat{\varphi }}`$ are constant, from Eq. (6), and Eqs. (6) and (6) respectively. But from Eq. (6.49) it follows that $`S_{\widehat{t}\widehat{\theta }}=0`$, or $`A=B=C=D=0`$, from Eq. (6.2), so that $`S_{\widehat{\theta }\widehat{\varphi }}=0`$ and $`S_{\widehat{r}\widehat{\theta }}=0`$ as well, from Eqs. (6.1) and (6) respectively. This contradicts the assumption $`S_{\widehat{r}\widehat{\theta }}0`$ so only the case $`S_{\widehat{r}\widehat{\theta }}=0`$ remains to be considered.
If $`S_{\widehat{r}\widehat{\theta }}=0`$, the set of equations (6.3) reduces to
$`0`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{r}}}{\mathrm{d}\tau _U}}\left[{\displaystyle \frac{m}{\nu S_{\widehat{t}\widehat{r}}S_{\widehat{r}\widehat{\varphi }}}}+\gamma \nu {\displaystyle \frac{\zeta _K}{\nu _K}}\right]S_{\widehat{t}\widehat{\varphi }},`$ (6.55)
$`0`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{\varphi }}}{\mathrm{d}\tau _U}}m{\displaystyle \frac{\nu S_{\widehat{r}\widehat{\varphi }}S_{\widehat{t}\widehat{r}}}{\nu S_{\widehat{t}\widehat{r}}S_{\widehat{r}\widehat{\varphi }}}}+\gamma {\displaystyle \frac{\zeta _K}{\nu _K}}[\nu S_{\widehat{t}\widehat{r}}\nu _K^2S_{\widehat{r}\widehat{\varphi }}],`$ (6.56)
provided that $`\nu S_{\widehat{t}\widehat{r}}S_{\widehat{r}\widehat{\varphi }}0`$. Substituting Eqs. (1.30), (6.7) and then Eqs. (A.3), (A.3) into Eq. (6.56) we obtain
$$\mathrm{\hspace{0.17em}0}=\zeta _K^2(\nu ^2\nu _K^2)[\nu _K^2S_{\widehat{t}\widehat{r}}^2S_{\widehat{r}\widehat{\varphi }}^2]+\nu _K^2c_0[\nu S_{\widehat{t}\widehat{r}}S_{\widehat{r}\widehat{\varphi }}]\gamma \nu _K\zeta _Kc_m(\nu ^2\nu _K^2)[S_{\widehat{t}\widehat{r}}\nu S_{\widehat{r}\widehat{\varphi }}].$$
(6.57)
Substituting Eq. (6) into Eq. (6.57) then gives
$$c_1=c_2=c_3=c_4=0,$$
(6.58)
implying
$$S_{\widehat{t}\widehat{\varphi }}=0,\frac{\mathrm{d}S_{\widehat{t}\widehat{r}}}{\mathrm{d}\tau _U}=0=\frac{\mathrm{d}S_{\widehat{r}\widehat{\varphi }}}{\mathrm{d}\tau _U},$$
(6.59)
from Eq. (6), and Eqs. (6) and (6) respectively. Hence, Eq. (6.55) is identically satisfied; moreover
$`c_0^{(\pm )}`$ $`=`$ $`{\displaystyle \frac{c_m}{2}}{\displaystyle \frac{\gamma \nu }{\gamma _K^2}}{\displaystyle \frac{\zeta _K}{\nu _K}}\{[2\nu _K^2(15\nu _K^2)](\nu ^2\nu _K^2)^2+3\nu _K^2\{\nu ^2[3\nu _K^2(7\nu _K^2)]`$
$`+\nu _K^4(4\nu _K^2)\}\pm {\displaystyle \frac{\nu _K}{\gamma ^2\nu }}[{\displaystyle \frac{\nu ^2}{\gamma _K^2}}+\nu _K^2(2+\nu _K^2)][\nu ^2(13\nu _K^2+4\nu ^2)8\nu _K^4]^{1/2}\}`$
$`\left\{\left[{\displaystyle \frac{\nu ^2}{\gamma _K^2}}+\nu _K^2(2+\nu _K^2)\right]^29\nu ^2\nu _K^4\right\}^1.`$ (6.60)
Next substituting Eqs. (6) and (6) and then Eq. (6.3) into Eq. (6.7) leads to
$`c_m^{(\pm )}`$ $`=`$ $`{\displaystyle \frac{m}{2}}{\displaystyle \frac{\{\nu _K^2\nu ^4\nu ^2[1\nu _K^2(3\nu _K^2)]\}^1}{\nu ^2+2\nu _K^2}}\{2{\displaystyle \frac{\nu ^6}{\gamma _K^2}}\nu ^4[1\nu _K^2(3+\nu _K^2)]`$
$`+\nu ^2\nu _K^2[2\nu _K^2(18\nu _K^2)]+4\nu _K^4(2+\nu _K^2)`$
$`\pm {\displaystyle \frac{\nu }{\nu _K}}\{\nu ^2[1\nu _K^2(3+\nu _K^2)]+\nu _K^2(2+\nu _K^4)\}[\nu ^2(13\nu _K^2+4\nu ^2)8\nu _K^4]^{1/2}\}.`$ (6.61)
Finally substituting Eqs. (6.3) and (6.3) into Eqs. (6) and (6), we obtain expressions for the only nonvanishing components of the spin tensor
$`S_{\widehat{t}\widehat{r}}`$ $`=`$ $`{\displaystyle \frac{m}{2\gamma }}{\displaystyle \frac{\nu _K}{\zeta _K}}{\displaystyle \frac{\nu ^2\nu _K^2}{\nu ^2+2\nu _K^2}}{\displaystyle \frac{\nu ^2(23\nu _K^2)+4\nu _K^2\pm \nu \nu _K[\nu ^2(13\nu _K^2+4\nu ^2)8\nu _K^4]^{1/2}}{\nu _K^2(\nu ^42)\nu ^2[1\nu _K^2(3\nu _K^2)]}},`$ (6.62)
$`S_{\widehat{r}\widehat{\varphi }}`$ $`=`$ $`{\displaystyle \frac{m}{2\gamma }}{\displaystyle \frac{\nu _K}{\zeta _K}}{\displaystyle \frac{\{\nu _K^2(\nu ^42)\nu ^2[1\nu _K^2(3\nu _K^2)]\}^1}{\nu ^2+2\nu _K^2}}\{\nu \nu _K[(\nu ^2+2\nu _K^2)(4\nu ^23\nu _K^2)`$
$`2(\nu ^2\nu _K^2)^2]\pm [\nu ^2(13\nu _K^2)2\nu _K^2][\nu ^2(13\nu _K^2+4\nu ^2)8\nu _K^4]^{1/2}\},`$ (6.63)
which are in agreement with the condition $`\nu S_{\widehat{t}\widehat{r}}S_{\widehat{r}\widehat{\varphi }}0`$ assumed above. This solution, having constant spin components, was already found in previous work .
Eq. (6.59) together with the fact that $`S_{\widehat{t}\widehat{\theta }}=0`$, $`S_{\widehat{r}\widehat{\theta }}=0`$ show that the total 4-momentum $`P`$ (see Eq. (3)) also lies in the cylinder of the circular orbit
$$P=mU+m_sE_{\widehat{\varphi }},$$
(6.64)
with
$$m_s=\gamma \frac{\zeta _K}{\nu _K}(\nu S_{\widehat{t}\widehat{r}}\nu _K^2S_{\widehat{r}\widehat{\varphi }}).$$
(6.65)
It can therefore be written in the form $`P=\mu U_p`$ with
$$U_p=\gamma _p[e_{\widehat{t}}+\nu _pe_{\widehat{\varphi }}],\nu _p=\frac{\nu +m_s/m}{1+\nu m_s/m},\mu =\frac{\gamma }{\gamma _p}(m+\nu m_s),$$
(6.66)
where $`\gamma _p=(1\nu _p^2)^{1/2}`$, provided that $`m+\nu m_s0`$. The T supplementary conditions can then be written as
$$S_{\widehat{t}\widehat{\varphi }}=0,S_{\widehat{r}\widehat{t}}+S_{\widehat{r}\widehat{\varphi }}\nu _p=0,S_{\widehat{\theta }\widehat{t}}+S_{\widehat{\theta }\widehat{\varphi }}\nu _p=0,$$
(6.67)
the last condition being identically satisfied, and with the equivalent azimuthal velocity $`\nu _p`$ given by
$$\nu _p^{(\pm )}=\frac{1}{2}\frac{\nu _K}{\nu ^2+2\nu _K^2}\{3\nu \nu _K\pm [\nu ^2(13\nu _K^2+4\nu ^2)8\nu _K^4]^{1/2}\},$$
(6.68)
from Eqs. (6.62) and (6.3). The reality condition of (6.68) requires that $`\nu `$ takes values outside the interval $`(\widehat{\nu },\widehat{\nu })`$, with $`\widehat{\nu }=\nu _K\sqrt{2}\sqrt{13+3\sqrt{33}}/40.727\nu _K`$; moreover, the timelike condition for $`|\nu _p|<1`$ is satisfied for all values of $`\nu `$ outside the same interval.
From (6.67) the spin vector orthogonal to $`U_p`$ is just $`\gamma _{p}^{}{}_{}{}^{1}S_{\widehat{r}\widehat{\varphi }}E_{\widehat{\theta }}`$. The spin-curvature force (3.7) turns out to be radially directed
$`F^{(\mathrm{sc})}=\gamma \zeta _K^2\left[2S_{\widehat{t}\widehat{r}}+\nu S_{\widehat{r}\widehat{\varphi }}\right]e_{\widehat{r}}.`$ (6.69)
The term on the left hand side of Eq. (1.1) can be written as
$$\frac{DP}{\mathrm{d}\tau _U}=[m\kappa m_s\tau _1]e_{\widehat{r}},$$
(6.70)
so that the balance equation (3.8) reduces to
$$ma(U)_{\widehat{r}}=F_{\widehat{r}}^{(\mathrm{so})}+F_{\widehat{r}}^{(\mathrm{sc})},$$
(6.71)
where
$$ma(U)_{\widehat{r}}=m\kappa ,F_{\widehat{r}}^{(\mathrm{so})}=m_s\tau _1,F_{\widehat{r}}^{(\mathrm{sc})}=\gamma \zeta _K^2\left[2S_{\widehat{t}\widehat{r}}+\nu S_{\widehat{r}\widehat{\varphi }}\right].$$
(6.72)
## 7 Conclusions
Spinning test particles in circular motion around a Schwarzschild black hole have been discussed in the framework of the Mathisson-Papapetrou approach supplemented by the usual standard conditions. One finds that apart from very special (and indeed artificially constrained) orbits where the spin tensor is closely matched to the curvature and torsion properties of the world line of the test particle or the static observer spin vector is constant and orthogonal to the plane of the orbit, the assumption of circular motion is not compatible with these equations. Indeed even in the former case, the test particle assumption is then violated except in the limit of massless particles following null geodesics, where the spin vector must be aligned with the direction of motion from general considerations. The spin-curvature force generically forces the motion away from circular orbits, so one needs a much more complicated machinery to attempt to study explicit solutions of this problem, solutions which must break the stationary axisymmetry.
## Appendix A Derivation of the solutions of the equations of motion
This Appendix derives the solutions of the equations of motion alone Eqs. (3.2)โ(3.4) and (3)โ(3) for the spin tensor along a circular orbit without supplementary conditions imposed. Standard elimination and differentiation techniques are used to find decoupled second order linear constant coefficient equations for certain spin components, from which one may calculate the remaining spin components that do not already satisfy decoupled first order such equations.
### A.1 The $`\nu =0`$ case
When the particle is at rest $`\nu =0`$ relative to the static observers, these equations reduce to
$`0`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}S_{\widehat{r}\widehat{\varphi }}}{\mathrm{d}\tau _U}}\nu _K\zeta _KS_{\widehat{t}\widehat{\varphi }},`$ (1.1)
$`0`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}S_{\widehat{\theta }\widehat{\varphi }}}{\mathrm{d}\tau _U}},`$ (1.2)
$`0`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}S_{\widehat{r}\widehat{\theta }}}{\mathrm{d}\tau _U}}\nu _K\zeta _KS_{\widehat{t}\widehat{\theta }},`$ (1.3)
$`0`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}m}{\mathrm{d}\tau _U}}+\nu _K\zeta _K{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{r}}}{\mathrm{d}\tau _U}},`$ (1.4)
$`0`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}^2S_{\widehat{t}\widehat{r}}}{\mathrm{d}\tau _U^2}}2\zeta _K^2S_{\widehat{t}\widehat{r}}+m\nu _K\zeta _K,`$ (1.5)
$`0`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}^2S_{\widehat{t}\widehat{\theta }}}{\mathrm{d}\tau _U^2}}+\zeta _K^2S_{\widehat{t}\widehat{\theta }}\zeta _K\nu _K{\displaystyle \frac{\mathrm{d}S_{\widehat{r}\widehat{\theta }}}{\mathrm{d}\tau _U}},`$ (1.6)
$`0`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}^2S_{\widehat{t}\widehat{\varphi }}}{\mathrm{d}\tau _U^2}}+\zeta _K^2S_{\widehat{t}\widehat{\varphi }}\zeta _K\nu _K{\displaystyle \frac{\mathrm{d}S_{\widehat{r}\widehat{\varphi }}}{\mathrm{d}\tau _U}}.`$ (1.7)
Eq. (1.2) implies that $`S_{\widehat{\theta }\widehat{\varphi }}=c_1`$ is constant. Solving Eq. (1.3) for $`\mathrm{d}S_{\widehat{r}\widehat{\theta }}/\mathrm{d}\tau _U`$, and substituting the result into Eq. (1.6) leads to the decoupled equation
$$0=\frac{\mathrm{d}^2S_{\widehat{t}\widehat{\theta }}}{\mathrm{d}\tau _U^2}+\omega _0^2S_{\widehat{t}\widehat{\theta }},\omega _0=\frac{\zeta _K}{\gamma _K}=\sqrt{\frac{M(r3M)}{r^3(r2M)}}.$$
(1.8)
Similarly solving Eq. (1.1) for $`\mathrm{d}S_{\widehat{r}\widehat{\varphi }}/\mathrm{d}\tau _U`$, and substituting the result into Eq. (1.7) leads to an analogous decoupled equation
$$0=\frac{\mathrm{d}^2S_{\widehat{t}\widehat{\varphi }}}{\mathrm{d}\tau _U^2}+\omega _0^2S_{\widehat{t}\widehat{\varphi }}.$$
(1.9)
Eq. (1.4) leads immediately to $`m=\nu _K\zeta _KS_{\widehat{t}\widehat{r}}+c_m`$ and substituting this into (1.5) yields
$$0=\frac{\mathrm{d}^2S_{\widehat{t}\widehat{r}}}{\mathrm{d}\tau _U^2}+\omega _1^2S_{\widehat{t}\widehat{r}}+\nu _K\zeta _Kc_m,\omega _1=\zeta _K(2+\nu _K^2)^{1/2}=\sqrt{\frac{M(2r3M)}{r^3(r2M)}}.$$
(1.10)
The three second order constant coefficient Eqs. (1.8), (1.9) and (1.10) are easily integrated, from which expressions for the remaining components of the spin tensor are then obtained from Eqs. (1.1) and (1.3). These have either oscillatory or exponential solutions depending on whether the squared frequencies in Eqs. (1.8) and (1.10) are positive or negative, or linear solutions when zero. This distinguishes the two intervals $`2M<r<3M`$ and $`r>3M`$, whose corresponding solutions are given by Eqs. (1) and (3) respectively. The special case $`r=3M`$ can be easily discussed by setting $`\nu _K=1`$ and $`\zeta _K=\sqrt{3}/(9M)`$ in Eqs. (1.1)โ(1.7). The corresponding solution is given by Eq. (2).
### A.2 The $`\nu =\pm \nu _K`$ case
When the particle moves along a geodesic with $`\nu =\pm \nu _K`$, Eqs. (3.2)โ(3.4) and (3)โ(3) simplify to
$`0`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}S_{\widehat{r}\widehat{\varphi }}}{\mathrm{d}\tau _U}}\nu _K{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{r}}}{\mathrm{d}\tau _U}},`$ (1.11)
$`0`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}S_{\widehat{\theta }\widehat{\varphi }}}{\mathrm{d}\tau _U}}\nu _K{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{\theta }}}{\mathrm{d}\tau _U}}{\displaystyle \frac{\zeta _K}{\gamma _K}}S_{\widehat{r}\widehat{\theta }},`$ (1.12)
$`0`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}S_{\widehat{r}\widehat{\theta }}}{\mathrm{d}\tau _U}}\pm \zeta _K\gamma _K[S_{\widehat{\theta }\widehat{\varphi }}\nu _KS_{\widehat{t}\widehat{\theta }}],`$ (1.13)
$`0`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}^2S_{\widehat{t}\widehat{\varphi }}}{\mathrm{d}\tau _U^2}}\pm {\displaystyle \frac{1}{\nu _K}}{\displaystyle \frac{\mathrm{d}m}{\mathrm{d}\tau _U}}\pm 2\zeta _K\gamma _K\left[{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{r}}}{\mathrm{d}\tau _U}}\nu _K{\displaystyle \frac{\mathrm{d}S_{\widehat{r}\widehat{\varphi }}}{\mathrm{d}\tau _U}}\right],`$ (1.14)
$`0`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}^2S_{\widehat{t}\widehat{r}}}{\mathrm{d}\tau _U^2}}\nu _K{\displaystyle \frac{\mathrm{d}^2S_{\widehat{r}\widehat{\varphi }}}{\mathrm{d}\tau _U^2}}2{\displaystyle \frac{\zeta _K}{\gamma _K}}{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{\varphi }}}{\mathrm{d}\tau _U}}3\zeta _K^2S_{\widehat{t}\widehat{r}},`$ (1.15)
$`0`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}^2S_{\widehat{t}\widehat{\theta }}}{\mathrm{d}\tau _U^2}}\nu _K{\displaystyle \frac{\mathrm{d}^2S_{\widehat{\theta }\widehat{\varphi }}}{\mathrm{d}\tau _U^2}}+\zeta _K^2[S_{\widehat{t}\widehat{\theta }}\pm 2\nu _KS_{\widehat{\theta }\widehat{\varphi }}],`$ (1.16)
$`0`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}^2S_{\widehat{t}\widehat{\varphi }}}{\mathrm{d}\tau _U^2}}\pm \nu _K{\displaystyle \frac{\mathrm{d}m}{\mathrm{d}\tau _U}}\pm 2\zeta _K\gamma _K\left[{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{r}}}{\mathrm{d}\tau _U}}\nu _K{\displaystyle \frac{\mathrm{d}S_{\widehat{r}\widehat{\varphi }}}{\mathrm{d}\tau _U}}\right].`$ (1.17)
The difference of Eqs. (1.14) and (1.17) leads to $`\mathrm{d}m/\mathrm{d}\tau _U=0`$, so $`m=c_m`$ is constant, which then implies from the same equations that
$$0=\frac{\mathrm{d}^2S_{\widehat{t}\widehat{\varphi }}}{\mathrm{d}\tau _U^2}\pm 2\zeta _K\gamma _K\left[\frac{\mathrm{d}S_{\widehat{t}\widehat{r}}}{\mathrm{d}\tau _U}\nu _K\frac{\mathrm{d}S_{\widehat{r}\widehat{\varphi }}}{\mathrm{d}\tau _U}\right].$$
(1.18)
Integration of Eq. (1.11) yields
$$S_{\widehat{r}\widehat{\varphi }}=\pm \nu _KS_{\widehat{t}\widehat{r}}+c_1,$$
(1.19)
and using this to replace $`S_{\widehat{r}\widehat{\varphi }}`$ in Eq. (1.18) and then integrating gives
$$\frac{\mathrm{d}S_{\widehat{t}\widehat{\varphi }}}{\mathrm{d}\tau _U}\pm 2\frac{\zeta _K}{\gamma _K}S_{\widehat{t}\widehat{r}}=c_2.$$
(1.20)
Using Eqs. (1.19) and (1.20) to replace $`S_{\widehat{r}\widehat{\varphi }}`$ and $`\mathrm{d}S_{\widehat{t}\widehat{\varphi }}/\mathrm{d}\tau _U`$ in Eq. (1.15) then leads to the decoupled second order equation
$$0=\frac{\mathrm{d}^2S_{\widehat{t}\widehat{r}}}{\mathrm{d}\tau _U^2}+\omega _2^2S_{\widehat{t}\widehat{r}}\pm 2\gamma _K\zeta _Kc_2,\omega _2=\zeta _K(43\gamma _K^2)^{1/2}=\sqrt{\frac{M(r6M)}{r^3(r3M)}}.$$
(1.21)
Taking the $`\tau _U`$ derivative of Eq. (1.12) and using it and Eq. (1.13) in Eq. (1.16) leads to the second order equation
$$0=\frac{\mathrm{d}^2S_{\widehat{t}\widehat{\theta }}}{\mathrm{d}\tau _U^2}+\omega _3^2S_{\widehat{t}\widehat{\theta }}+3\zeta _K^2\gamma _K^2\nu _KS_{\widehat{\theta }\widehat{\varphi }},\omega _3=\zeta _K=\left(\frac{M}{r^3}\right)^{1/2}.$$
(1.22)
Finally using Eq. (1.22) and Eq. (1.13) in the equation obtained by taking the $`\tau _U`$ derivative of Eq. (1.12) gives a second decoupled second order equation
$$0=\frac{\mathrm{d}^2S_{\widehat{\theta }\widehat{\varphi }}}{\mathrm{d}\tau _U^2}+\omega _4^2S_{\widehat{\theta }\widehat{\varphi }},\omega _4=\zeta _K(3\gamma _K^22)^{1/2}=\sqrt{\frac{M}{r^3(r3M)}}.$$
(1.23)
Integrating the two decoupled second order equations (1.21) and (1.23), one can then integrate Eq. (1.22) too. The remaining components of the spin tensor are then determined by Eqs. (1.13), (1.19), and (1.20). Note that the frequencies $`\omega _3`$ and $`\omega _4`$ agree only for $`\gamma _K=1`$, or $`\nu _K=0`$, which would imply $`M=0`$: so this special case is not relevant.
Now the two intervals $`3M<r<6M`$ and $`r>6M`$ have differing signs for the squared angular velocities, and the corresponding solutions are given by Eqs. (1) and (3) respectively. The special case $`r=6M`$ can be easily handled as well, setting $`\nu _K=1/2`$ and $`\zeta _K=\sqrt{6}/(36M)`$ in Eqs. (1.11)โ(1.17). The corresponding solution is given by Eq. (2).
### A.3 The general case: $`\nu 0`$ and $`\nu \pm \nu _K`$
Here we deal with the general form of Eqs. (3.2)โ(3.4) and (3)โ(3). Solving Eqs. (3.2)โ(3.4) for their first terms and substituting these derivative terms into Eqs. (3)โ(3), one finds
$`{\displaystyle \frac{\mathrm{d}^2S_{\widehat{t}\widehat{\varphi }}}{\mathrm{d}\tau _U^2}}`$ $`=`$ $`{\displaystyle \frac{\gamma }{\nu }}{\displaystyle \frac{\zeta _K}{\nu _K}}[(2\nu _K^21)\nu ^2\nu _K^2]{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{r}}}{\mathrm{d}\tau _U}}\gamma ^2\zeta _K^2(\nu ^2\nu _K^2)S_{\widehat{t}\widehat{\varphi }}{\displaystyle \frac{1}{\nu }}{\displaystyle \frac{\mathrm{d}m}{\mathrm{d}\tau _U}},`$ (1.24)
$`{\displaystyle \frac{\mathrm{d}^2S_{\widehat{t}\widehat{r}}}{\mathrm{d}\tau _U^2}}`$ $`=`$ $`\gamma ^3\nu {\displaystyle \frac{\zeta _K}{\nu _K}}[\nu ^2+\nu _K^22]{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{\varphi }}}{\mathrm{d}\tau _U}}\gamma ^4{\displaystyle \frac{\zeta _K^2}{\nu _K^2}}[(3\nu _K^21)\nu ^22\nu _K^2]S_{\widehat{t}\widehat{r}}`$
$`\gamma ^4\nu \zeta _K^2(\nu ^2\nu _K^2)S_{\widehat{r}\widehat{\varphi }}+m\gamma ^3{\displaystyle \frac{\zeta _K}{\nu _K}}(\nu ^2\nu _K^2),`$ (1.25)
$`{\displaystyle \frac{\mathrm{d}^2S_{\widehat{t}\widehat{\theta }}}{\mathrm{d}\tau _U^2}}`$ $`=`$ $`{\displaystyle \frac{\gamma ^3}{\gamma _K^2}}{\displaystyle \frac{\zeta _K}{\nu _K}}\nu ^2{\displaystyle \frac{\mathrm{d}S_{\widehat{r}\widehat{\theta }}}{\mathrm{d}\tau _U}}{\displaystyle \frac{\gamma ^4}{\gamma _K^2}}\zeta _K^2S_{\widehat{t}\widehat{\theta }}+\gamma ^4\nu {\displaystyle \frac{\zeta _K^2}{\nu _K^2}}[(1+2\nu _K^2)\nu ^23\nu _K^2]S_{\widehat{\theta }\widehat{\varphi }},`$ (1.26)
$`{\displaystyle \frac{\mathrm{d}^2S_{\widehat{t}\widehat{\varphi }}}{\mathrm{d}\tau _U^2}}`$ $`=`$ $`\gamma \nu {\displaystyle \frac{\zeta _K}{\nu _K}}[\nu ^2+\nu _K^22]{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{r}}}{\mathrm{d}\tau _U}}\gamma ^2{\displaystyle \frac{\zeta _K^2}{\nu _K^2}}(\nu ^2\nu _K^2)[\nu ^2+\nu _K^21]S_{\widehat{t}\widehat{\varphi }}\nu {\displaystyle \frac{\mathrm{d}m}{\mathrm{d}\tau _U}}.`$ (1.27)
Solving Eqs. (1.24) and (1.27) for $`\mathrm{d}m/\mathrm{d}\tau _U`$ and $`\mathrm{d}^2S_{\widehat{r}\widehat{\varphi }}/\mathrm{d}\tau _U^2`$, one finds
$`{\displaystyle \frac{\mathrm{d}m}{\mathrm{d}\tau _U}}`$ $`=`$ $`\gamma {\displaystyle \frac{\zeta _K}{\nu _K}}(\nu ^2\nu _K^2)\left[{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{r}}}{\mathrm{d}\tau _U}}\gamma \nu {\displaystyle \frac{\zeta _K}{\nu _K}}S_{\widehat{t}\widehat{\varphi }}\right],`$ (1.28)
$`{\displaystyle \frac{\mathrm{d}^2S_{\widehat{t}\widehat{\varphi }}}{\mathrm{d}\tau _U^2}}`$ $`=`$ $`2{\displaystyle \frac{\gamma \nu }{\gamma _K^2}}{\displaystyle \frac{\zeta _K}{\nu _K}}{\displaystyle \frac{\mathrm{d}S_{\widehat{t}\widehat{r}}}{\mathrm{d}\tau _U}}+{\displaystyle \frac{\gamma ^2}{\gamma _K^2}}{\displaystyle \frac{\zeta _K^2}{\nu _K^2}}(\nu ^2\nu _K^2)S_{\widehat{t}\widehat{\varphi }}.`$ (1.29)
Solving Eq. (3.2) for $`S_{\widehat{t}\widehat{\varphi }}`$, leads to
$$S_{\widehat{t}\widehat{\varphi }}=\frac{1}{\gamma }\frac{\nu _K}{\zeta _K}\frac{1}{\nu ^2\nu _K^2}\left[\nu \frac{\mathrm{d}S_{\widehat{t}\widehat{r}}}{\mathrm{d}\tau _U}\frac{\mathrm{d}S_{\widehat{r}\widehat{\varphi }}}{\mathrm{d}\tau _U}\right].$$
(1.30)
Using Eq. (1.30) in Eq. (1.28) leads to
$$m=\gamma \frac{\zeta _K}{\nu _K}[\nu S_{\widehat{r}\widehat{\varphi }}\nu _K^2S_{\widehat{t}\widehat{r}}]+c_m.$$
(1.31)
Then using Eqs. (1.30) and (1.31) in Eq. (A.3) leads to
$`0`$ $`=`$ $`[(12\nu _K^2)\nu ^2+\nu _K^2]{\displaystyle \frac{\mathrm{d}^2S_{\widehat{t}\widehat{r}}}{\mathrm{d}\tau _U^2}}+\nu [\nu ^2+\nu _K^22]{\displaystyle \frac{\mathrm{d}^2S_{\widehat{r}\widehat{\varphi }}}{\mathrm{d}\tau _U^2}}+{\displaystyle \frac{\gamma ^2}{\gamma _K^2}}\nu {\displaystyle \frac{\zeta _K^2}{\nu _K^2}}(\nu ^2\nu _K^2)^2S_{\widehat{r}\widehat{\varphi }}`$
$`+\gamma ^2{\displaystyle \frac{\zeta _K^2}{\nu _K^2}}(\nu ^2\nu _K^2)[(14\nu _K^2)\nu ^2+\nu _K^2(2+\nu _K^2)]S_{\widehat{t}\widehat{r}}+\gamma {\displaystyle \frac{\zeta _K}{\nu _K}}(\nu ^2\nu _K^2)^2c_m.`$ (1.32)
Using Eq. (1.30) in Eq. (1.29) leads to
$$\frac{\mathrm{d}^2S_{\widehat{r}\widehat{\varphi }}}{\mathrm{d}\tau _U^2}\nu \frac{\mathrm{d}^2S_{\widehat{t}\widehat{r}}}{\mathrm{d}\tau _U^2}\frac{\gamma ^2}{\gamma _K^2}\frac{\zeta _K^2}{\nu _K^2}(\nu ^2\nu _K^2)[S_{\widehat{r}\widehat{\varphi }}+\nu S_{\widehat{t}\widehat{r}}]=c_0.$$
(1.33)
Then solving these last two equations for $`\mathrm{d}^2S_{\widehat{r}\widehat{\varphi }}/\mathrm{d}\tau _U^2`$ and $`\mathrm{d}^2S_{\widehat{t}\widehat{r}}/\mathrm{d}\tau _U^2`$ gives
$`{\displaystyle \frac{\mathrm{d}^2S_{\widehat{t}\widehat{r}}}{\mathrm{d}\tau _U^2}}`$ $`=`$ $`\gamma ^2{\displaystyle \frac{\zeta _K^2}{\nu _K^2}}\left[{\displaystyle \frac{\nu ^2}{\gamma _K^2}}\nu _K^2(2+\nu _K^2)\right]S_{\widehat{t}\widehat{r}}2{\displaystyle \frac{\gamma ^2}{\gamma _K^2}}\nu {\displaystyle \frac{\zeta _K^2}{\nu _K^2}}S_{\widehat{r}\widehat{\varphi }}`$
$`+\gamma ^2\nu {\displaystyle \frac{\nu ^2+\nu _K^22}{\nu ^2\nu _K^2}}c_0+\gamma ^3{\displaystyle \frac{\zeta _K}{\nu _K}}(\nu ^2\nu _K^2)c_m,`$ (1.34)
$`{\displaystyle \frac{\mathrm{d}^2S_{\widehat{r}\widehat{\varphi }}}{\mathrm{d}\tau _U^2}}`$ $`=`$ $`{\displaystyle \frac{\gamma ^2}{\gamma _K^2}}{\displaystyle \frac{\zeta _K^2}{\nu _K^2}}(\nu ^2+\nu _K^2)S_{\widehat{r}\widehat{\varphi }}+\gamma ^2\nu \zeta _K^2(1+2\nu _K^2)S_{\widehat{t}\widehat{r}}`$
$`\gamma ^2{\displaystyle \frac{\nu ^2(12\nu _K^2)+\nu _K^2}{\nu ^2\nu _K^2}}c_0+\gamma ^3\nu {\displaystyle \frac{\zeta _K}{\nu _K}}(\nu ^2\nu _K^2)c_m,`$ (1.35)
whose solution is given by Eqs. (6) and (6). Then substituting these solutions into Eq. (1.30), one obtains (6).
Substituting Eqs. (3.3) and (3.4) into Eq. (3.11) gives
$$0=\frac{\mathrm{d}^2S_{\widehat{t}\widehat{\theta }}}{\mathrm{d}\tau _U^2}+\mathrm{\Omega }_1^2S_{\widehat{t}\widehat{\theta }}+3\gamma ^2\nu \zeta _K^2S_{\widehat{\theta }\widehat{\varphi }},\mathrm{\Omega }_1^2=\gamma ^2\frac{\zeta _K^2}{\gamma _K^2}.$$
(1.36)
Taking the $`\tau _U`$ derivative of Eq. (3.3) and using Eqs. (1.36) and (3.4), one gets
$$0=\frac{\mathrm{d}^2S_{\widehat{\theta }\widehat{\varphi }}}{\mathrm{d}\tau _U^2}+\mathrm{\Omega }^2S_{\widehat{\theta }\widehat{\varphi }},\mathrm{\Omega }^2=\gamma ^2\nu ^2\frac{\zeta _K^2}{\nu _K^2}(1+2\nu _K^2).$$
(1.37)
This is easily integrated to give the solution (6.1). The remaining components of the spin can be obtained directly from Eqs. (1.36) and (3.4) yielding Eqs. (6.2) and (6). The frequencies $`\mathrm{\Omega }`$ and $`\mathrm{\Omega }_1`$ agree for the particular value of the azimuthal velocity given by
$$\nu _0=\pm \frac{\nu _K}{\gamma _K}(1+2\nu _K^2)^{1/2}.$$
(1.38)
The solution corresponding to this special case is given by Eqs. (6.12)โ(6) together with Eqs. (6)โ(6.7) evaluated at $`\nu =\nu _0`$.
## References
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# Astrometric Microlensing Constraints on a Massive Body in the Outer Solar System with Gaia
## 1 Introduction
Observations of long-period comets in the inner Solar System suggest not only a substantial population of comets at 50,000 to 100,000 AU (the Oort Cloud; Oort 1950), but a mechanism for effectively perturbing the orbits of these comets. Such a perturber must be massive enough to hold considerable gravitational influence on the Oort cloud. Galactic tidal perturbations could be the cause of a steady steam of cometary infall (Byl, 1983) while close encounters with passing stars would cause a more punctuated cascade (Hills, 1981). Punctuated (and perhaps periodic; Hut et al. 1987) cometary showers into the inner Solar system could also be caused by a perturber that is bound to the Sun. Specific predictions of the mass and orbit ($`0.003M_{}`$, $`d110\times 10^4`$ AU) of such a perturber depend on whether its existence is invoked to explain temporal features in mass extinctions on Earth (โNemesisโ; e.g., Davis et al. 1984,Whitmire & Jackson 1984, and Vandervoort & Sather 1993) and/or the trajectories of anomalous streams of comets (โPlanet Xโ; see Murray 1999 and Matese et al. 1999, but see a more cautious view from Horner & Evans 2002).
There are some direct constraints on the existence of any massive (planetary or larger) perturber in the outer Solar System. To have eluded detection by all-sky synoptic surveys like Hipparcos and Tycho-2 (Hรธg et al., 2000), any massive body in the outer Solar System but must be fainter than $`V11`$ mag, corresponding to absolute magnitude $`M_V>21`$ mag for $`d<0.1`$ pc. This constraint rules out main sequence stars above the hydrogen-burning limit.
Detection of a massive perturber through reflected Solar light grows increasingly difficult with increasing distance due to $`r^4`$ dimming. In reflected light, at current sensitivity limits and angular size coverages, discoveries of objects in the Kuiper Belt at $``$40 AU have only recently become routine (e.g., Brown et al. 2004). Yet even with an all-sky synoptic survey to limiting magnitudes of $`R=24`$ mag (e.g., Pan-STARRS<sup>1</sup><sup>1</sup>1http://pan-starrs.ifa.hawaii.edu/project/reviews/PreCoDR/documents/ scienceproposals/sol.pdf), massive planets like Neptune would be undetectable via reflected light beyond $``$800 AU and a 0.1 $`M_{}`$ perturber with a density of 1 g cm<sup>-3</sup> would be undetected with $`d>2000`$ AU.
Old and cooled degenerate stars (emitting thermally) could be faint enough to have gone undetected. The oldest neutron star (NS) known with an apparent thermal emission component is B0950$`+`$08 with $`M_B(20.0\pm 0.2)`$ mag (Zharikov et al., 2002) ($`d260\mathrm{pc}`$; age = 10<sup>7.2</sup> yr). At $`d=90,000`$ AU, the source would be $`B13`$ mag, likely detectable with the next generation synoptic surveys. However, with a cooling time that of the age of the Solar System, we would expect a NS perturber to have cooled considerably, likely to $`T<10^3`$ K from $`T10^5`$ K (extrapolating from Page et al. 2004) and so would be significantly fainter than current detection levels. Constraints on the existence of even colder distant planetary-mass objects from the lack of detection of their thermal infrared emission with the Infrared Astronomical Satellite (IRAS) are largely superseded by constraints from the ephemerides of the outer planets (Hogg et al., 1991). An infrared survey with significantly higher spatial resolution and sensitivity may provide interesting constraints on distant objects.
Surveys that monitor distant stars with high cadence to search for occultations by foreground objects are in principle sensitive to objects of mass as low as $`0.01M_{}`$ out to the Galactic tidal radius of the solar system at $`10^5\mathrm{AU}`$. However, the probability that any one object will occult a sufficiently bright background star to be detectable is very low. Therefore, in order to detect any occultation events at all, a large number of objects must be present. Thus such surveys can only constrain the existence of a substantial population of objects, and will place essentially no constraints on the existence of individual bodies in the outer solar system.
Clearly, the limits on faint massive objects in the outer Solar System must be probed with a fundamentally different technique than through reflected, thermally emitted, or occulted light. Here we suggest an indirect search for massive outer Solar System bodies by observing the differential astrometric microlensing signature that such bodies would impart on the distant stars. As the apparent position of the lens moves on the sky, astrometric monitoring of background sources in the vicinity of the lens (with the appropriate sensitivity) will reveal a complex pattern of apparent motion of those background sources. In ยง2 we introduce the microlensing formalism in the regime of interest. Detecting the astrometric microlensing signature of a lens requires either the background stars to move and/or the lens to move. Nearby objects exhibit extremely large parallaxes and so the apparent position of the lens, regardless of whether it can be detected directly in reflected light, sweeps out a large area of influence on the sky even if the proper motion of lens is small. Indeed, parallax dominates the apparent motion of objects in Solar orbit beyond the Kuiper Belt. In ยง3 we estimate the detectability of a nearby massive perturber using the data from the Gaia mission<sup>2</sup><sup>2</sup>2Launch excepted June 2011; http://astro.estec.esa.nl/GAIA/ using a Monte Carlo simulation. Finally, we highlight some improvements in the detectability estimate for future work.
## 2 Properties of Induced Parallax
Consider a distant source with parallax $`\mathrm{\Pi }_S`$ with an (angular) separation $`\theta `$ from a foreground massive body with parallax $`\mathrm{\Pi }_X`$. The foreground body will deflect the apparent position of the centroid of the background source relative to its unlensed position by,
$$\mathrm{\Delta }\theta =\frac{๐ฎ}{u^2+2}\theta _\mathrm{E},$$
(1)
where $`๐ฎ=\theta /\theta _\mathrm{E}`$ is the angular separation of lens and source in units of the angular Einstein ring radius,
$$\theta _\mathrm{E}=(\kappa M_X\mathrm{\Pi }_{rel})^{1/2}.$$
(2)
Here $`\kappa =4G/c^2\mathrm{AU}=8.144\mathrm{mas}/M_{}`$, and $`\mathrm{\Pi }_{rel}=\mathrm{\Pi }_X\mathrm{\Pi }_S`$ is the relative lens-source parallax. For the cases considered here, $`\mathrm{\Pi }_S\mathrm{\Pi }_X`$. For $`u1`$, $`|\mathrm{\Delta }\theta |=\theta _\mathrm{E}^2/\theta `$.
Due to parallax, the apparent position of the massive body will trace out an ellipse on the sky over the course of a year. In addition, it will have a proper motion $`\mu _X`$ due to its intrinsic motion. In ecliptic coordinates, the position of the lens at time $`t`$, relative to time $`t_0`$ has components,
$$\delta \lambda _X(t)=\mathrm{\Pi }_X\mathrm{sin}(\omega [tt_0])+\mu _X(tt_0)\mathrm{cos}\gamma $$
(3)
$$\delta \beta _X(t)=\mathrm{\Pi }_X\mathrm{sin}(\beta )\mathrm{cos}(\omega [tt_0])+\mu _X(tt_0)\mathrm{sin}\gamma ,$$
(4)
where $`\beta `$ is the ecliptic latitude of the object, and $`\gamma `$ is the angle of the proper motion with respect to the ecliptic plane. For orbits with zero inclination (in the plane of the ecliptic), $`\gamma =0`$. We have also assumed $`\mathrm{\Pi }_X1\mathrm{rad}`$.
The deflection tracks of background stars that are astrometrically microlensed by the motion of lens parallax (hereafter โinduced parallaxโ) can exhibit a variety of shapes depending on the angular position with respect to the parallactic ellipse of the lens. Figure 1 shows a realization of several tracks around a neutron star at 10,000 AU. For sources at large impact parameter to the lens, the apparent positions over the year trace out a curved path along a distortion angle approximately parallel with the direction of motion of the lens at the minimum impact of the source along the parallactic ellipse. Near the position of maximal parallactic position of the lens, these curves resemble โtear dropโ shapes. For impacts comparable to the semi-minor axis of the parallactic ellipse ($`\mathrm{\Pi }_X\mathrm{sin}\beta `$), the deflection tracks take the appearance of โcrescentโ shapes or a โcircle-within-circleโ. This is due to comparable deflection during the nearest impact and the distant opposite-side impact months later; these such types of deflection paths are obviously more common at smaller $`|\beta |`$. Sources interior to the parallactic ellipse are deflected near maximally twice a year, resulting in shapes resembling a โfigure eightโ.
Although we call the deflection tracks due to parallactic motion of the lens โinduced parallax,โ the deflection tracks generally do not resemble the traditional parallactic ellipse. First, the eccentricity of the tracks does not generally scale with $`\mathrm{cos}b`$. Second, the direction of motion along the tracks is retrograde with respect to the parallactic motion of the lens. Moreover, unlike traditional parallax (where the date of maximum departure is fixed by the ecliptic azimuth), the time of maximum departure from the unlensed positions depends only on the time of minimum impact of the source to the lens. In these ways, the source parallactic motion may be distinguished from the effects due to induced parallax in principle. However, in practice the presence of intrinsic source proper motion and parallax, which are typically much larger than the signals we are concerned with here, as well as poor sampling and signal-to-noise ratio, may cause considerable degradation of the detectability. We consider these issues in more detail below.
## 3 Estimating the Lens Mass-Distance Sensitivity of an Astrometric Survey
Figure 1 shows a rather dramatic effect of a nearby neutron star upon a background field, with deflections of many background sources more than arcseconds from unlensed positions. Since the magnitude of the deflection tracks scales as the mass of the lens, all-sky astrometric missions could, in principle, probe to masses significantly smaller than $`M_{}`$. We now quantify what mass/distance configurations would give rise to a detectable signal in the presence of astrometric uncertainty and a finite number of position samples of the background sources. Though the deflection of a single background source may not be detectable, clearly neighboring sources will exhibit similar, correlated deflection; therefore, the presence of a nearby massive lens can be inferred at a statistically significant level by aggregating a collection of statistically insignificant deflections.
Consider a massive body in solar orbit with mass $`M_X`$ and heliocentric distance $`D_X`$. This body will have a parallax $`\mathrm{\Pi }_X=\mathrm{AU}/D_X`$, and a proper motion $`\mu _X=v_X/D_X`$, where $`v_X`$ is its transverse velocity. If we assume that the body is in a circular orbit, and that $`D_X\mathrm{AU}`$ (so that projection effects are small), then $`v_X=v_{}\mathrm{\Pi }_X^{1/2}`$.
Now consider that the body is moving in front of a background screen of source stars with surface density $`\mathrm{\Sigma }_{}`$, and that series of $`N`$ astrometric measurements of these stars are taken at times $`t_j`$. At each time $`t_j`$, we can compute the deflection due to the lens $`\mathrm{\Delta }\theta _k(t_j)=[\mathrm{\Delta }\theta _{\lambda ,k}(t_j),\mathrm{\Delta }\theta _{\beta ,k}(t_j)]`$, for each source $`k`$, using the expressions presented in ยง2. Assuming all the source stars have the same (one-dimensional) astrometric uncertainty $`\sigma _{ast}`$, we can estimate the total signal-to-noise ratio $`\mathrm{S}/\mathrm{N}`$ with which the deflection of the massive body is detected as,
$$(\mathrm{S}/\mathrm{N})^2=\frac{1}{2\sigma _{ast}^2}\underset{k}{}\underset{j}{}(\mathrm{\Delta }\theta _{\lambda ,k}(t_j)\mathrm{\Delta }\theta _{\lambda ,k})^2+(\mathrm{\Delta }\theta _{\beta ,k}(t_j)\mathrm{\Delta }\theta _{\beta ,k})^2.$$
(5)
Here $`\mathrm{\Delta }\theta _{\lambda ,k}`$ and $`\mathrm{\Delta }\theta _{\lambda ,k}`$ are the average deflections, i.e. $`\mathrm{\Delta }\theta _{\lambda ,k}N^1_j\theta _{\lambda ,j}`$. These are the average positions of the source determined over the course of the Gaia mission relative to some external reference grid well away from the deflector. Adopting this $`\mathrm{S}/\mathrm{N}`$ criterion for detection is in some sense conservative, in that it only defines the significance with which the positions of the background stars differ from the null hypothesis of no deflections. The effective $`\mathrm{S}/\mathrm{N}`$ will likely be increased by fitting a model to the data which implicitly accounts for the shape of the deflection track, as well as the correlation between neighboring sources. We note that, using the median deflections in equation (5), rather than the mean, increases the $`\mathrm{S}/\mathrm{N}`$ by $`10\%`$.
We estimate the $`\mathrm{S}/\mathrm{N}`$ using a simple Monte Carlo<sup>3</sup><sup>3</sup>3Note that it is possible, using some simplifying assumptions and by analyzing the problem in limiting regimes, to make significant analytical progress and arrive at simple expressions for the signal-to-noise ratio as a function of the mass and distance to the perturber, as well as the surface density and astrometric accuracy of the source stars. We have chosen not to present these analytic expressions here, as there are not fully general, and so one ultimately must resort to numerical evaluations to determine the detectability in all relevant regimes. We note that these analytic results generally confirm the numerical results we now present.. We create a random screen of stars, and simulate a series of $`N`$ uniformly sampled measurements. We then calculate $`\mathrm{S}/\mathrm{N}`$ using equation (5). Under our assumptions, the $`\mathrm{S}/\mathrm{N}`$ depends on the parameters of the lens, $`M_X,\mathrm{\Pi }_X,\beta ,t_0,\gamma `$, as well as the properties of the source stars, $`\mathrm{\Sigma }_{},\sigma _{ast}`$. We calculate $`\mathrm{S}/\mathrm{N}`$ for many different realizations of the positions of the background source stars, and we also vary the input parameters. We find the following approximate expression for the signal-to-noise ratio,
$$\mathrm{S}/\mathrm{N}\{\begin{array}{cc}\frac{10\mu \mathrm{as}}{\sqrt{2}\sigma _{ast}}\left(\frac{M_X}{M_{}}\right)\left(\frac{D_X}{10^3\mathrm{AU}}\right)^1\left(\frac{\mathrm{\Sigma }_{}}{10^3\mathrm{arcsec}^2}\right)^{1/2}\left(\frac{N}{40}\right)^{1/2}\left(1+\mathrm{sin}\beta \right)\hfill & \mathrm{if}\mathrm{\Sigma }_{}\pi \mathrm{\Pi }_X^21\hfill \\ \frac{10\mu \mathrm{as}}{\sqrt{2}\sigma _{ast}}\left(\mathrm{\Sigma }_{}\pi \mathrm{\Pi }_X^2\right)^{1/2}\left(\frac{M_X}{M_{}}\right)\left(\frac{D_X}{10^3\mathrm{AU}}\right)^1\left(\frac{\mathrm{\Sigma }_{}}{10^3\mathrm{arcsec}^2}\right)^{1/2}\left(\frac{N}{40}\right)^{1/2}\left(1+\mathrm{sin}\beta \right)\hfill & \mathrm{if}\mathrm{\Sigma }_{}\pi \mathrm{\Pi }_X^2<1\hfill \end{array}.$$
(6)
The two regimes in equation (6) correspond to the strong, โcollisionalโ regime where there is on average one star in the parallactic ellipse, and the weak, โtidalโ regime where there is typically less than one star in the ellipse. Equation 6 is generally accurate to considerably better than the variance at fixed values of the parameters due to Poisson fluctuations in the number density and location of source stars, for most parameter combinations. The $`\mathrm{S}/\mathrm{N}`$ can vary by a large amount due to Poisson noise depending on the parameters, and especially so in the tidal regime for low $`\mathrm{\Sigma }_{}`$. Note that, as reflected in equation (6), we find that the $`\mathrm{S}/\mathrm{N}`$ does not depend on $`t_0`$ or $`\gamma `$ to within the Poisson fluctuations.
### 3.1 Application to Gaia
In order to provide a quantitative estimate of the mass-distance sensitivity of an astrometric survey to massive objects in the outer solar system, we adopt parameters appropriate for the Gaia mission. Gaia will monitor the entire sky synoptically for five years, acquiring astrometric measurements for $`O(10^9)`$ stars down to apparent magnitudes of $`V20`$. For bright stars $`(V12)`$, Gaia will have a single-measurement astrometric precision limit of $`30\mu \mathrm{as}`$, whereas at $`V20`$, the astrometric accuracy will be $`1400\mu \mathrm{as}`$. Typically, each star will have $`100200`$ astrometric measurements, grouped in clusters of 2 to 5 measurements each.
To proceed with our estimate, we adopt a model of the surface density of source stars on the sky as a function of magnitude, Galactic latitude and longitude, and a model of the expected astrometric performance of Gaia. This allows us to predict the total $`\mathrm{S}/\mathrm{N}`$ with which a object of a given mass and distance would be detected with Gaia, at a given location in the sky.
The expected performance of Gaia has and will continue to evolve, and the final mission astrometric accuracy is therefore impossible to access currently. For definiteness, we assume that the (one-dimensional) astrometric uncertainty of each measurement is given by,
$$\sigma _{1D}^2(V)=\{\begin{array}{cc}\sigma _{sys}^2\hfill & \mathrm{if}V12.5\hfill \\ \sigma _s^210^{0.4(V12.5)}+\sigma _b^210^{0.8(V20)}\hfill & \mathrm{if}V>12.5\hfill \end{array},$$
(7)
with $`\sigma _{sys}=\sigma _s=30\mu \mathrm{as}`$ and $`\sigma _b=1000\mu \mathrm{as}`$. This form was chosen to reproduce the astrometric accuracies from Table 1 of Belokurov & Evans (2002). Gaia will not make astrometric measurements uniformly across the sky; certain ecliptic latitudes will be sampled a larger number of times than others. We assume that the number of samples as a function of ecliptic latitude $`\beta `$ is given by,
$$N_{samp}=100+300\mathrm{exp}\left[\left(\left|\frac{|\beta |35^{}}{10^{}}\right|\right)^{1/2}\right].$$
(8)
This form was chosen to qualitatively reproduce Figure 5 of Belokurov & Evans (2002). We assume that the samples are clustered into groups of $`n_c`$ points, and so the effective number of points is $`N=N_{samp}/n_c`$, and the effective astrometric accuracy of each point is $`\sigma _{ast}=\sigma _{1D}(V)/\sqrt{n_c}`$. This assumes that the single-measurement errors can be reduced by root-$`n`$ averaging. This may not be the case: the measurement errors in any given cluster may be correlated, or there may exist systematic errors that are not reducible. Since it is difficult to anticipate the behavior of the astrometric errors in advance, we will adopt the assumption of root-$`n`$ averaging for simplicity. We adopt $`n_c=5`$ (Belokurov & Evans, 2002). For other values of $`n_c`$, the $`\mathrm{S}/\mathrm{N}`$ for any given star, as well as the integrated $`\mathrm{S}/\mathrm{N}`$, will scale as $`\sqrt{n_c/5}`$.
We determine the surface density of source stars as a function of position and magnitude using a simple model for the Galaxy. For the density distribution of sources, we adopt the double-exponential disk plus barred bulge model of Han & Gould (1995, 2003). We assume that the dust column is independent of Galactocentric radius and has an exponential distribution in height above the plane with a scale height of $`120\mathrm{pc}`$. We normalize the midplane column density so that the $`V`$-band extinction is $`A_V=1\mathrm{mag}(D_s/\mathrm{kpc})`$, where $`D_s`$ is the distance to the source. We also will show results assuming the dust model of Belokurov & Evans (2002), which is similar to ours for $`\beta 20^{}`$, but differs in detail for latitudes closer to the plane. Finally, we assume a $`V`$-band luminosity function that is independent position and is equal to the Bahcall-Soneira (Bahcall & Soneira, 1980) luminosity function for $`M_V10`$, and is constant for $`10M_V20`$.
The surface density of stars down to $`V=20`$ in our model ranges from $`10^5\mathrm{arcsec}^2`$ near the Galactic poles, to $`10^3\mathrm{arcsec}^2`$ near the Galactic anticenter, to a maximum of $`0.1\mathrm{arcsec}^2`$ within a few degrees of the Galactic center. Therefore, regions of the sky near the Galactic plane and especially the Galactic center will have greater sensitivity to lower mass and/or more distant perturbers for fixed $`\mathrm{S}/\mathrm{N}`$. The total number of stars in the sky with $`10V20`$ in this model is $`1.3\times 10^9`$ for our standard dust extinction model, and $`1.0\times 10^9`$ for the Belokurov & Evans (2002) dust model. Thus the average surface density is $`10^3\mathrm{arcsec}^2`$.
Figure 2 shows the distribution of $`\mathrm{S}/\mathrm{N}`$ for an object with $`M=3000M_{}10M_{\mathrm{jup}}`$ and $`D=10^4\mathrm{AU}`$ located in three different locations on the sky: near the Galactic bulge, anticenter and north Galactic pole. The source densities in these three locations vary considerably, from $`\mathrm{\Sigma }_{}10^2\mathrm{arcsec}^2`$ near the Galactic bulge to $`10^5\mathrm{arcsec}^2`$ near the pole. The shape of the distribution of $`\mathrm{S}/\mathrm{N}`$ depends on the location on the sky, through the distribution of source densities as a function of magnitude (and so astrometric accuracy). For the locations near the Galactic plane with high source densities, the distribution of $`\mathrm{S}/\mathrm{N}`$ has a tail toward higher values, and so the total $`\mathrm{S}/\mathrm{N}`$ is generally dominated by one or two stars. For the location near the Galactic pole, a larger number of stars contribute significantly to the total $`\mathrm{S}/\mathrm{N}`$. The total $`\mathrm{S}/\mathrm{N}`$ (integrated over $`V`$-magnitude from $`V=10`$ to $`V=20`$) for these three locations are $`(\mathrm{S}/\mathrm{N})_{\mathrm{tot}}=94.4`$ (bulge), 16.5 (anticenter), and 1.4 (pole).
Figure 3 shows contours of constant $`(\mathrm{S}/\mathrm{N})_{\mathrm{tot}}`$ for a object with $`M=3000M_{}10M_{\mathrm{jup}}`$ and $`D=10^4\mathrm{AU}`$. The distribution of $`(\mathrm{S}/\mathrm{N})_{\mathrm{tot}}`$ on the sky is highly non-uniform: objects of a given $`M`$ and $`D_X`$ located toward certain regions of the sky will be detected with higher $`(\mathrm{S}/\mathrm{N})_{\mathrm{tot}}`$ than if they are located in other regions. The $`(\mathrm{S}/\mathrm{N})_{\mathrm{tot}}`$ is primarily driven by the surface density of stars, and therefore regions of the sky near the Galactic plane and especially the Galactic center are preferred. However, it is also the case that the number of samples $`N_{samp}`$ depends on ecliptic latitude, such that stars with ecliptic latitude $`\pm 35^{}`$ will have several times more astrometric measurements than stars near the ecliptic poles. Therefore locations near ecliptic latitudes of $`\pm 35^{}`$ will also have higher $`(\mathrm{S}/\mathrm{N})_{\mathrm{tot}}`$ for fixed perturber mass and distance.
Figure 4 shows the fraction of the sky enclosed by contours of a given $`(\mathrm{S}/\mathrm{N})_{\mathrm{tot}}`$, i.e. the fraction of the sky over which an object of mass $`M=3000M_{}`$ and distance $`D_X=10^4\mathrm{AU}`$ would be detected with $`\mathrm{S}/\mathrm{N}`$ greater than a given value. We determine the fraction of sky above a given $`(\mathrm{S}/\mathrm{N})_{\mathrm{tot}}`$ for a range of masses and distances. Objects with mass greater than a minimum mass
$$M_{min}(290,490,750)M_{}\left\{\begin{array}{cc}\left(\frac{D_X}{10^4\mathrm{AU}}\right)\left[\frac{(\mathrm{S}/\mathrm{N})_{\mathrm{th}}}{5}\right]\hfill & \mathrm{if}D_XD_b\hfill \\ \left(\frac{D_X}{D_b}\right)\left(\frac{D_X}{10^4\mathrm{AU}}\right)\left[\frac{(\mathrm{S}/\mathrm{N})_{\mathrm{th}}}{5}\right]\hfill & \mathrm{if}D_X>D_b\hfill \end{array}\right\},\mathrm{for}f_{\mathrm{sky}}=(10\%,50\%,100\%).$$
(9)
can be detected with $`\mathrm{S}/\mathrm{N}(\mathrm{S}/\mathrm{N})_{\mathrm{th}}`$, where $`f_{\mathrm{sky}}`$ is the fraction of the sky. Here $`D_b`$ is the โbreak distanceโ, and has values of $`D_b=(4470,1550,780)\mathrm{AU}`$ for $`f_{\mathrm{sky}}=(10\%,50\%,100\%)`$. These limits are shown in Figure 5.
Figure 4 also shows the fraction of the sky within $`10^{}`$ of the ecliptic plane enclosed by contours of a given $`(\mathrm{S}/\mathrm{N})_{\mathrm{tot}}`$. Since the ecliptic plane fortuitously passes near the Galactic bulge, the slope for this curve is shallower than that for the entire sky, resulting in a relatively larger fraction of the area for which a high $`\mathrm{S}/\mathrm{N}`$ detections are possible.
There are several obvious limitations of our calculations. One is that we have neglected the motion of background stars due to parallax. To the extent that these motions correlate with the microlensing signal, they will tend to degrade the signal-to-noise ratio, by effectively allowing one to partially โfit outโ the anomalous excursions. Motions of stars in binaries could also confound a clean measurement of induced parallax. Also, since we adopted the simple scaling relation in equation (6) when integrating over the magnitude distribution of source stars, we have neglected the effect of the Poisson fluctuations of the surface density and location of stars on the total signal-to-noise ratio. To provide a rough estimate of the magnitude of these effects, we have performed a few simulations where we determine the signal-to-noise ratio for stars of a given magnitude directly from the Monte Carlo simulation (which per force includes Poisson fluctuations), while explicitly fitting for the parallax of the source stars. Since these calculations are extremely time intensive, we have not performed a comprehensive exploration, but rather checked only a few cases. For these few cases, we find that fitting for the parallax of the source does indeed reduce $`(\mathrm{S}/\mathrm{N})_{\mathrm{tot}}`$, but by a relatively small factor, $`10\%`$. On the other hand, we find that the effect of Poisson fluctuations causes us to underestimate $`(\mathrm{S}/\mathrm{N})_{\mathrm{tot}}`$, by as much as $`75\%`$.
## 4 Discussion and Conclusions
We have shown that a substantial, as yet unexplored, region of mass-distance parameter space of nearby massive bodies will be accessible with the current incarnation of the datastream from the Gaia experiment. We have focused on the effect of โinduced parallaxโ caused only by the parallax of the lens as it sweeps through the parallactic ellipse. Based on our albeit simplistic simulation, the search for massive bodies in the outer Solar System by the observation of induced parallax has a reasonable chance of uncovering the proposed perturber of cometary orbits in the Oort cloud (Figs. 4 and 5). In particular, we believe that the non-detection of a massive body in the Gaia dataset using the proposed technique would relegate the proposed mass-distances of Planet X to a significantly smaller parameter space then the currently allowed space<sup>4</sup><sup>4</sup>4It is noteworthy that Horner & Evans (2002) also appeal to Gaia for constraining the existence of Planet X, but by making use of ephemeris data of $``$1000 long-period comets that would be discovered by Gaia relatively uniformly over the sky..
Murray (1999) made specific predictions for the current position of โPlanet Xโ on the sky, based on the clustering of cometary aphelion distances. Since the $`\mathrm{S}/\mathrm{N}`$ map of the sky is not uniform, it is interesting to ask with what $`\mathrm{S}/\mathrm{N}`$ one would expect to detect โPlanet Xโ with the allowed mass and distances, at its expected position. Figure 3 shows the positional error ellipse from Murray (1999). The expected $`\mathrm{S}/\mathrm{N}`$ for $`M_X=10^3M_{}`$ and $`D_X=10^4\mathrm{AU}`$ ranges from $`(\mathrm{S}/\mathrm{N})_{\mathrm{tot}}3`$ to $`(\mathrm{S}/\mathrm{N})_{\mathrm{tot}}12`$. The mass/distance limit for thresholds of $`(\mathrm{S}/\mathrm{N})_{th}=1,3`$, and $`5`$ in this error ellipse is shown in Figure 6; roughly 25% the allowed parameter space could be excluded at $`3\sigma `$ with a non-detection.
Hypothesis for the mass ($`0.03M_{}`$) and distance ($`10^5\mathrm{AU}`$) of Nemesis will likely be difficult to test with Gaia (see Fig. 6), due primarily to the large distance and thus small size of the parallactic ellipse. However, specific predictions for the current position of Nemesis might be testable using a targeted astrometric satellite with higher astrometric precision than Gaia, such as the Space Interferometry Mission (SIM). Of course, constraints on smaller-mass objects at any distance could be obtained with with an all-sky synoptic experiment that has improved astrometric accuracy but with a similar limiting magnitude (โSuperGaiaโ, Fig. 6) or by probing more stars to fainter levels with Gaia-like astrometric accuracies.
Should a significant detection be made, what can be learned about the lens? In principle, the astrometric data alone provide an estimate of the mass, position, distance, and proper motion of the lens for high-$`\mathrm{S}/\mathrm{N}`$ detections of induced parallax for stars very near to the parallactic ellipse. Orbit determination will generally be difficult, unless there is a significant acceleration over the five year mission lifetime; this is only expected for relatively nearby lenses. For more modest $`\mathrm{S}/\mathrm{N}`$ detections, or detections in the tidal regime where the source stars are quite distant from the parallactic ellipse, the information will be seriously degraded, and degeneracies between the mass, distance, and angular separation from the lens arise. In the extreme case where only one distant star is significantly perturbed, the detection may yield very little information about the lens. Exploration of the information that can be extracted from these various classes of detections is beyond the scope of this paper, but is an interesting topic for future study.
Further follow-up of potential candidates may be possible with a variety of methods. Astrometric follow-up of individual background sources may be possible with SIM with higher astrometric precision and cadence than possible with Gaia; such measurements may improve on the determination of the lens parameters. Direct detection of the reflected light from some candidates may be possible with ultra-deep imaging using very large aperture, next generation, ground-based, optical/near-infrared telescopes such as the Giant Magellan Telescope (GMT), the Thirty Meter Telescope (TMT), or the Overwhelmingly Large Telescope (OWL). Finally, the James Webb Space Telescope (JWST) should have the sensitivity to detect the thermal emission from essentially all objects detected astrometrically by Gaia.
A similar astrometric microlensing search with Gaia for massive stellar remnants in the Solar neighborhood ($`d150`$ pc) was proposed by Belokurov & Evans (2002) but with several important differences compared to the present work. First, we considered the detectability of an object significantly closer to Earth so that the lens parallax is $`10^{57}`$ larger than the typical source parallaxes whereas that difference is only $`10^{12}`$ for Solar neighborhood lenses. We also focused on Solar System lenses in Solar orbit where the parallax motion dominates proper motion; the motion of Solar neighborhood objects are dominated by proper motion. Both these different regimes result in significantly different microlensing tracks of a single background star (compare our Fig. 1 with Fig. 1 of Belokurov & Evans 2002). Second, we focus on detection of objects with a planet-scale mass whereas the analysis technique of Belokurov & Evans 2002 is optimized to constrain the mass function of stellar-mass objects in the Solar neighborhood (see, e.g., Fig. 3) with $`M>0.1M_{}`$. Last, and conceptually the most distinct, we consider the detectability of a single massive object using the aggregate induced parallax signatures of thousands of stars whereas Belokurov & Evans focused on constraining the properties of a large population of faint stellar-mass objects, where the mass of each object is inferred using the astrometric microlensing โeventโ a single background source. Ultimately, though, both analyzes make use of the same datastream and act toward complimentary goals.
We have assumed that our lenses are point-like, and so have ignored the effects of occultation of the background sources by the lens. If the angular size of the lens $`\theta _X`$ is an appreciable fraction of its angular Einstein ring radius $`\theta _\mathrm{E}`$, then both occultation and lensing effects can potentially be important (Agol, 2002; Takahashi, 2003). In Figure 5, we show the locus of mass and distance where $`\theta _X=\theta _\mathrm{E}`$. Objects with $`M_X10^3M_{}`$ will have angular sizes that are larger than their Einstein ring radii provided they are closer than $`4000\mathrm{AU}`$; for such objects, complete occultations are possible. However, an occultation will obviously only occur if a background source happens to be located within an angular radius of the lens when a measurement is taken. This condition is met when the number of measurements satisfies $`N\mathrm{\Sigma }_{}\pi \theta _X^21`$. Figure 5 shows the region of parameter space for which at least one measurement will be occulted by the lens, for typical background source densities of $`\mathrm{\Sigma }_{}=2\times 10^3\mathrm{arcsec}^2`$. Clearly, for most lenses, occultation effects are negligible.
In our simulation, we assumed the perturber is in a circular orbit around the Sun. However, we found that our results are essentially independent of the proper motion of the lens. Furthermore, realistic motions along the line-of-sight are unlikely to alter our signal-to-noise ratio estimates substantially for the distances considered herein. Therefore, the assumption that the lens is on a circular orbit or indeed even bound to the Sun is immaterial to our conclusions.
As we have discussed, an obvious shortcoming of our estimation is that we have neglected the motion of background stars due to parallax, proper motion, and orbits. These motions will tend to degrade the signal-to-noise ratio, effectively introducing more free parameters to help explain away anomalous excursions. Still our preliminary calculations show that source parallax is not likely to degrade the $`\mathrm{S}/\mathrm{N}`$ substantially, however these calculations were admittedly limited. We hope to perform a more comprehensive study to quantify the effect of a realistic background screen in future work.
Our simplistic simulation for S/N estimation also neglects another feature of data that could be exploited to improve the S/N. Any nearby foreground massive source will lens multiple source background stars differently in the course of a 5 year mission. Moreover, neighboring background sources will be lensed similarly. So the expectation of correlated deflection paths (which are fixed for a given lens mass, distance, and proper motion) could be used to create a โmatched filterโ for improving the sensitivity of detecting a nearby massive lens. Though computationally very expensive, one can envision applying such a filter to the Gaia dataset for all possible nearby lens masses at all possible distances and positions on sky to search for a signal. Aside from the need to simultaneously constrain the parallax, proper motion, and orbital parameters of all background sources, the matched filter search may also need to search for a possible changing parallax of the lens over the mission lifetime: a massive object passing nearby that is unbound to the Sun with $`|v|30`$ km s<sup>-1</sup> would travel $``$30 AU over 5 yr, with some of this motion in the radial direction from the Sun.
Finally, the choice of the appropriate $`\mathrm{S}/\mathrm{N}`$ threshold for a robust detection deserves some discussion. Here one must not only consider the astrometric noise properties of the sources, but also the total number of independent trials performed in searching the data with a matched filter. This latter factor can be quite crucial in the current context, given the fact that one is performing a blind search over the entire sky with $`O(10^9)`$ source stars, with many independent filters corresponding to varying lens locations, distances, masses, and proper motions.
While a high signal-to-noise ratio measurement of the entire induced parallax of single star will yield the lens mass, sky position, proper motion, and distance, the likelihood of such a special configuration is rare. Instead, each of these lens events will contribute individually to constraints on the lens properties at different times, leading to the possibility of improving the signal-to-noise ratio of the lens properties through the matched filter. Another utility of global astrometric filtering of the Gaia data is that the masses and ephemerides of known Solar System objects might be determined a priori, based solely on measurements of the astrometric microlensed background; whether the masses determined thusly will be more precise than measured by other means remains to the be seen.
BSG supported by a Menzel Fellowship from the Harvard College Observatory. JSB was partially supported by a grant from the Harvard-Smithsonian Center for Astrophysics. We would like to thank Avi Loeb for comments on the manuscript, and the anonymous referee for a prompt and helpful report. We would like to extend special thanks to Andy Gould for his detailed and comprehensive comments and insightful discussions which led to a much improved paper. Lastly, JSB thanks Eugene Chiang and Josh Eisner for their enthusiasm during the early stages of this work.
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# Asymptotic analysis of spatially inhomogeneous stiff and ultra-stiff cosmologies
## I Introduction
The conjectures of Belinskiว, Khalatnikov and Lifshitz (BKL) BKL assert that the structure of space-time singularities in general relativity (GR) have the following properties (for a perfect fluid with $`\gamma 2`$, where $`(\gamma 1)p/\rho `$ is defined as the ratio of the pressure $`p`$ to the energy density $`\rho `$): 1. Each spatial point evolves towards the singularity as if it were a spatially homogeneous cosmology. 2. Space-times with non-stiff matter, $`\gamma <2`$, have the property that asymptotically close to the singularity matter is not dynamically significant and the singularities in generic four dimensional space-times are space-like and oscillatory (mixmaster behavior). 3. In the case of stiff matter, $`\gamma =2`$, the matter is not insignificant near the singularity (leading to non-oscillatory behavior) and generically have anisotropic singularities which are space-like and non-oscillatory.
We shall investigate the BKL conjectures in cosmological models with an effective equation of state $`\gamma =2`$ and extend them to models with $`\gamma >2`$ by calculating the past asymptotic decay rates in general ($`G_0`$) spatially inhomogeneous models. These results are then supported by a numerical analysis of the behavior of spatially inhomogeneous solutions to Einsteinโs equations near an initial singularity in a special class of Abelian $`G_2`$ spatially inhomogeneous models.
There are a number of cosmological models of current physical interest which have an effective equation of state $`\gamma 2`$. Although a complete fundamental theory is not presently known the phenomenological consequences can be understood by studying an effective low-energy theory, which leads to the introduction of additional fields (e.g., scalar fields) in the high curvature regime close the Planck time scale. Scalar fields are believed to be abundant and pervasive in all fundamental theories of physics applicable in the early Universe, particularly in dimensionally reduced higher-dimensional theories Green1987 ; Olive1990a . In addition, scalar field cosmological models are of great importance in the study of the early Universe, particularly in the investigation of inflation Olive1990a ; inf and โquintessenceโ scalar field models Caldwell1998a consistent with observations of the present accelerated cosmic expansion PR . Superstring theory represents the most promising candidate for a unified theory of the fundamental interactions, including gravity Green1987 . It is widely believed that eleven-dimensional supergravity represents the low-energy limit of $`M`$-theory Witten95 . In the low-energy limit, to lowest-order in both the string coupling and the inverse string tension, all massive modes in the superstring spectrum decouple and only the massless sectors remain, which are determined by the corresponding supergravity actions. A definitive prediction of string theory is the existence of a scalar field, known as the dilaton, interpreted as a modulus field parametrizing the radius of the eleventh dimension. There are two further massless excitations that are common to all five perturbative string theories, namely the metric tensor field (the graviton) and an anti-symmetric form field. When the higher-dimensional metric is compactified, additional form fields and scalar moduli fields are produced. Thus the low-energy effective action of the theory essentially reduces to GR plus massless scalar fields. A massless scalar field (or moduli field etc.) close to the initial singularity has an effective equation of state $`\gamma =2`$.
Models in which $`\gamma 2`$ arise naturally in ekpyrotic Kho01A and cyclic Ste02 cosmological models, which have a big crunch/big bang transition with a contraction phase dominated by a scalar field with $`\gamma 2`$ to the future design . In particular, it was shown erik that if $`\gamma >2`$, chaotic mixmaster oscillations due to anisotropy and curvature are suppressed and the contraction is described by a spatially homogeneous and isotropic evolution. This result was subsequently generalized to theories where the scalar field couples to $`p`$-forms, and it was also shown that $`_2`$ orbifold compactification also contributes to suppressing chaotic behavior. Indeed, it was concluded that chaos is avoided in contracting heterotic $`M`$-theory models if $`\gamma >2`$ at the crunch.
There is currently great interest in higher-dimensional gravity theories inspired by string theory in which the matter fields are confined to a $`3+1`$-dimensional โbrane-worldโ embedded in higher dimensions, while the gravitational field can also propagate in the extra bulk dimensions rubakov . There has been particular interest in the dynamics of the Universe at early times in Randall-Sundrum-type brane-world cosmological models randall . A unique feature of brane cosmology is that $`\rho ^2`$ dominates at early times, leading to an effective equation of state parameter with $`\gamma >2`$, which will give rise to completely different behavior to that in GR. The cosmological implications of brane world models have been extensively investigated Maartens . In particular, it was found that an isotropic singularity is a past-attractor in all orthogonal Bianchi models Coley . Moreover, the asymptotic dynamical evolution of spatially inhomogeneous brane-world cosmological models close to the initial singularity was studied numerically CHL and it was found that there always exists an initial singularity, characterized by the fact that spatial derivatives are dynamically negligible, which is isotropic for all physical parameter values. The numerical results were supported by a qualitative dynamical analysis and a calculation of the past asymptotic decay rates CHL .
Andersson and Rendall Andersson2001 proved that a generic inhomogeneous cosmology with a stiff fluid or a scalar field tends to a velocity-dominated solution, which is the Jacobs solution (WainwrightHsu89, , p. 1426), along individual timelines. For a generic inhomogeneous cosmology with $`\gamma >2`$, we will show that it tends to a flat isotropic solution<sup>1</sup><sup>1</sup>1Namely the Binรฉtruy, Deffayet and Langlois solution BDL , which is essentially the flat Friedmann-Lemaรฎtre solution (WainwrightEllis97, , Ch. 2) with $`\gamma >2`$. along individual timelines. Note that these solutions are self-similar Carr . In this paper we shall show that these results are supported by the calculations of past asymptotic decay rates in general $`G_0`$ models (with $`\gamma 2`$)<sup>2</sup><sup>2</sup>2We assume that a massless scalar field (or a massive scalar field close to the initial singularity) can be modelled by a perfect fluid with a stiff equation of state. This is justified in detail in Andersson2001 . and numerical simulations in a class of $`G_2`$ models with one tilt degree of freedom.
## II Asymptotic dynamics at early times for $`G_2`$ cosmologies
We first consider the dynamics of a class of $`G_2`$ spatially inhomogeneous cosmological models with one spatial degree of freedom. The governing system of evolution equations constitute a system of autonomous partial differential equations in two independent variables. We follow the formalism of elst which utilizes area expansion-normalized scale-invariant dependent variables, and we use the separable area gauge to consider analytically and numerically the asymptotic evolution of the class of $`G_2`$ cosmologies with one tilt degree of freedom<sup>3</sup><sup>3</sup>3This class is described by the area expansion rate $`\beta H\frac{1}{2}\sigma _{11}`$ (thesis, , eq. (2.103)) and the normalized variables $`(E_1{}_{}{}^{1},A,N_\times ,N_{},\mathrm{\Sigma }_\times ,\mathrm{\Sigma }_2,\mathrm{\Sigma }_+,\mathrm{\Sigma }_{},\mathrm{\Omega },v)`$ defined with respect to a time-like congruence that corresponds to the separable area gauge. $`E_1^1`$ is the frame coefficient, $`A,N_\times ,N_{}`$ are the components of spatial curvatures, $`\mathrm{\Sigma }_\times ,\mathrm{\Sigma }_2,\mathrm{\Sigma }_+,\mathrm{\Sigma }_{}`$ are the shears, $`\mathrm{\Omega }`$ is the normalized density parameter, and $`v`$ is the tilt of the perfect fluid. See (thesis, , Appendix D) for the governing equations, where we set $`\mathrm{\Lambda }=0`$ and evolve $`\beta `$ instead (thesis, , eq. (2.108)). We remind the reader that the logarithmic time variable $`t\mathrm{ln}\mathrm{}`$ is used. near the cosmological initial singularity. For recent works on these models (with $`\gamma <2`$), we refer to thesis ; for vacuum $`G_2`$ models, we refer to Andersson05 and references therein. The decay rates at early times can be derived by following the analyses in Lim04 and CHL by exploiting asymptotic silence.<sup>4</sup><sup>4</sup>4See Uggla03 for the notion of the silent boundary and its role in past asymptotic dynamics. The results for $`G_2`$ models with one tilt degree of freedom are given below.
### II.1 Case $`\gamma =2`$
For the case $`\gamma =2`$, we assume that the following conditions hold uniformly for open intervals of $`x`$.
$`C_1:`$ $`\underset{t\mathrm{}}{lim}(E_1{}_{}{}^{1},A,N_\times ,N_{},\mathrm{\Sigma }_\times ,\mathrm{\Sigma }_2,v)=\mathrm{๐},\underset{t\mathrm{}}{lim}(\mathrm{\Sigma }_+,\mathrm{\Sigma }_{})=(\widehat{\mathrm{\Sigma }}_+(x),\widehat{\mathrm{\Sigma }}_{}(x))`$
$`C_2:`$ $`_x(E_1{}_{}{}^{1},A,N_\times ,N_{},\mathrm{\Sigma }_\times ,\mathrm{\Sigma }_2,v,\mathrm{\Sigma }_+,\mathrm{\Sigma }_{})\text{are bounded as}t\mathrm{}.`$
$`C_3:`$ $`V=๐ช(f(t))\text{implies}_xV=๐ช(f(t))\text{(asymptotic expansions in time}`$
$`\text{can be differentiated with respect to the spatial coordinates)}.`$
The decay rates as $`t\mathrm{}`$ are then given by
$$(E_1{}_{}{}^{1},A)=e^{2t}[(\widehat{E}_1{}_{}{}^{1},\widehat{A})+๐ช(f)],(\mathrm{\Sigma }_+,\mathrm{\Sigma }_{},\mathrm{\Omega })=(\widehat{\mathrm{\Sigma }}_+,\widehat{\mathrm{\Sigma }}_{},\widehat{\mathrm{\Omega }})+๐ช(g)$$
(1)
$$\mathrm{\Sigma }_2=e^{k_1t}[\widehat{\mathrm{\Sigma }}_2+๐ช(g)],\mathrm{\Sigma }_\times =e^{k_2t}[\widehat{\mathrm{\Sigma }}_\times +๐ช(h)],N_{}=e^{k_3t}[\widehat{N}_{}+๐ช(h)]$$
(2)
$$N_\times =e^{2t}[\widehat{E}_1{}_{}{}^{1}_{x}^{}\widehat{\mathrm{\Sigma }}_{}t+\widehat{N}_\times +๐ช(h)],v=e^{2t}[\frac{1}{2}\widehat{E}_1{}_{}{}^{1}_{x}^{}\mathrm{ln}\widehat{\mathrm{\Omega }}t+\widehat{v}+๐ช(h)]$$
(3)
where $`\widehat{\mathrm{\Sigma }}_+:=\frac{1}{2}[1\widehat{\mathrm{\Sigma }}_{}^2\widehat{\mathrm{\Omega }}]`$, and
$$k_1=3\widehat{\mathrm{\Sigma }}_++\sqrt{3}\widehat{\mathrm{\Sigma }}_{},k_2=2\sqrt{3}\widehat{\mathrm{\Sigma }}_{},k_3=2(1+\sqrt{3}\widehat{\mathrm{\Sigma }}_{}),$$
(4)
$$f=e^{4t}+e^{2k_1t},g=te^{4t}+e^{2k_1t}+e^{2k_2t}+e^{2k_3t},h=t(e^{4t}+e^{2k_1t}+e^{2k_2t}+e^{2k_3t}).$$
(5)
The area expansion rate $`\beta `$ and its associated deceleration parameter satisfy
$$\beta =e^{3(1\widehat{\mathrm{\Sigma }}_+)t}[\widehat{\beta }+๐ช(g)],q=2+๐ช(f).$$
(6)
The ten hat variables above are functions of $`x`$, and satisfy the two constraints
$$0=\widehat{E}_1{}_{}{}^{1}_{x}^{}\mathrm{ln}\widehat{\beta }+\widehat{r},0=\widehat{E}_1{}_{}{}^{1}_{x}^{}\mathrm{ln}\widehat{\mathrm{\Sigma }}_2\left[\widehat{r}+3\widehat{A}\sqrt{3}\widehat{N}_\times \right],$$
(7)
where $`\widehat{r}=3\widehat{A}\widehat{\mathrm{\Sigma }}_+3\widehat{N}_\times \widehat{\mathrm{\Sigma }}_{}+3\widehat{N}_{}\widehat{\mathrm{\Sigma }}_\times 3\widehat{v}\widehat{\mathrm{\Omega }}`$. That leaves eight hat variables, the same number as when we specify the initial conditions for numerical simulations.<sup>5</sup><sup>5</sup>5Only six of the hat variables are essential, since we can use the remaining temporal gauge freedom to set $`\widehat{A}=0`$, and parameterize $`x`$ to set $`\widehat{E}_1{}_{}{}^{1}=1`$. To ensure the convergence of $`N_{}`$, $`\mathrm{\Sigma }_\times `$ and $`\mathrm{\Sigma }_2`$, the exponents $`k_1`$, $`k_2`$ and $`k_3`$ must be positive for all $`x`$. This implies that the attractor is confined within the region
$$\mathrm{\Sigma }_{}^H>\frac{1}{\sqrt{3}}(1+\mathrm{\Sigma }_+^H),\mathrm{\Sigma }_{}^H<0,\mathrm{\Sigma }_{}^H>\sqrt{3}\mathrm{\Sigma }_+^H$$
(8)
inside the Kasner circle (triangle I in Figure 3a) in the state space of Hubble-normalized variables, where
$$\mathrm{\Sigma }_\pm ^H=\frac{\mathrm{\Sigma }_\pm }{1\mathrm{\Sigma }_+}.$$
(9)
Note that this restriction is gauge-dependent.
The area expansion-normalized Weyl scalar $`๐ฒ`$ is given by
$$๐ฒ^2=\frac{1}{9}(\widehat{\mathrm{\Sigma }}_+\widehat{\mathrm{\Sigma }}_{}^2)^2+\frac{1}{9}(13\widehat{\mathrm{\Sigma }}_+)^2\widehat{\mathrm{\Sigma }}_{}^2+๐ช(g+t^2e^{4t}).$$
(10)
Figure 1 shows the hat variables as computed at $`t=50`$ and $`t=100`$ for a numerical run (e.g., we plot $`e^{2t}E_1^1`$, $`e^{2t}(N_\times +tE_1{}_{}{}^{1}_{x}^{}\mathrm{\Sigma }_{})`$, etc).<sup>6</sup><sup>6</sup>6For numerical simulations, we use CLAWPACK. Since it handles exponential growth rather inaccurately, we evolve the following variables instead:
$$\mathrm{ln}E_1{}_{}{}^{1},A/E_1{}_{}{}^{1},\mathrm{\Sigma }_{},e^{2t}N_\times ,\mathrm{ln}\mathrm{\Sigma }_\times ,\mathrm{ln}N_{},\mathrm{ln}\mathrm{\Omega },e^{2t}\mathrm{artanh}(v),\mathrm{ln}\mathrm{\Sigma }_2,\mathrm{ln}\beta .$$
Compare with (thesis, , Appendix D). One should ensure that $`\mathrm{\Sigma }_\times `$ and $`N_{}`$ are positive during the simulation. We use 32768 grid points with periodic boundary condition, and run the simulation from $`t=0`$ to $`t=100`$. For the initial condition, we set $`\gamma =2`$, $`\beta _0=1`$ and
$$ฯต=0.1,(E_1{}_{}{}^{1})_0=1,A_0=0,(\mathrm{\Sigma }_{})_0=0.2,(\mathrm{\Sigma }_\times )_0=0.2,(N_{})_0=0.2,\mathrm{\Omega }_0=1,(\mathrm{\Sigma }_2)_0=0.4$$
in conjunction with (thesis, , eq. (9.1)). Since the plots at both times coincide, the decay rates are confirmed. Although the constraints tend to zero, they do not tend to zero fast enough and equation (7) is not satisfied to a sufficient degree of accuracy (this numerical problem can be solved by standard methods). Nonetheless, the numerical results are good enough for confirming the decay rates. The numerical simulations consequently confirm the calculations of the decay rates, provide evidence for the conditions $`C_1`$$`C_3`$, and hence lend support to the conjectures formulated above.
### II.2 Case $`\gamma >2`$
For the case $`\gamma >2`$, we assume that the following conditions hold uniformly for open intervals of $`x`$.
$`C_1:`$ $`\underset{t\mathrm{}}{lim}(E_1{}_{}{}^{1},A,N_\times ,N_{},\mathrm{\Sigma }_+,\mathrm{\Sigma }_{},\mathrm{\Sigma }_\times ,\mathrm{\Sigma }_2,v)=\mathrm{๐},`$
$`C_2:`$ $`_x(E_1{}_{}{}^{1},A,N_\times ,N_{},\mathrm{\Sigma }_\times ,\mathrm{\Sigma }_2,v,\mathrm{\Sigma }_+,\mathrm{\Sigma }_{})\text{are bounded as}t\mathrm{}.`$
$`C_3:`$ $`V=๐ช(f(t))\text{implies}_xV=๐ช(f(t))\text{(asymptotic expansions in time}`$
$`\text{can be differentiated with respect to the spatial coordinates)}.`$
The decay rates as $`t\mathrm{}`$ are then given by
$$(E_1{}_{}{}^{1},A,N_{},N_\times ,v)=\eta [(\widehat{E}_1{}_{}{}^{1},\widehat{A},\widehat{N}_{},\widehat{N}_\times ,\widehat{v})+๐ช(\xi )]$$
(11)
$$(\mathrm{\Sigma }_+,\mathrm{\Sigma }_{},\mathrm{\Sigma }_\times ,\mathrm{\Sigma }_2,\mathrm{\Omega }1)=\xi \left[(\widehat{\mathrm{\Sigma }}_+,\widehat{\mathrm{\Sigma }}_{},\widehat{\mathrm{\Sigma }}_\times ,\widehat{\mathrm{\Sigma }}_2,2\widehat{\mathrm{\Sigma }}_+)+๐ช(\xi )\right]$$
(12)
where
$$\eta =e^{\frac{1}{2}(3\gamma 2)t},\xi =e^{\frac{3}{2}(\gamma 2)t}.$$
(13)
The area expansion rate $`\beta `$ and its associated deceleration parameter satisfy
$$\beta =e^{\frac{3}{2}\gamma t}[\widehat{\beta }+๐ช(\xi )],q=\frac{1}{2}(3\gamma 2)+๐ช(\xi ).$$
(14)
The ten hat variables above are functions of $`x`$, and satisfy the two constraints
$$0=\widehat{E}_1{}_{}{}^{1}_{x}^{}\mathrm{ln}\beta \frac{3}{2}\gamma \widehat{v},0=\widehat{E}_1{}_{}{}^{1}_{x}^{}\mathrm{ln}\widehat{\mathrm{\Sigma }}_2(\frac{3}{2}\gamma \widehat{v}+3\widehat{A}\sqrt{3}\widehat{N}_\times ).$$
(15)
That leaves eight hat variables, the same number as when we specify the initial conditions for numerical simulations.
The area expansion-normalized Weyl scalar $`๐ฒ`$ is given by
$$๐ฒ^2=\frac{1}{9}(\widehat{\mathrm{\Sigma }}_+^2+\widehat{\mathrm{\Sigma }}_{}^2+\widehat{\mathrm{\Sigma }}_\times ^2+\widehat{\mathrm{\Sigma }}_2^2)\xi ^2+๐ช(\xi ^3).$$
(16)
Figure 2 shows the hat variables as computed at $`t=50`$ and $`t=100`$ for a numerical run (e.g. plot $`E_1{}_{}{}^{1}/\eta `$, etc).<sup>7</sup><sup>7</sup>7For the case $`\gamma >2`$, we evolve the variables
$$\mathrm{ln}E_1{}_{}{}^{1},A/E_1{}_{}{}^{1},\mathrm{\Sigma }_{}/\xi ,\mathrm{\Sigma }_\times /\xi ,N_{}/\eta ,N_\times /\eta ,(\mathrm{\Omega }1)/\xi ,\mathrm{artanh}(v)/\eta ,\mathrm{ln}\mathrm{\Sigma }_2,\mathrm{ln}\beta .$$
We use the same initial condition as before, except with $`\gamma =2.1`$. That the plots at both times coincide confirms the decay rates.<sup>8</sup><sup>8</sup>8We again note that numerically the constraints do not tend to zero sufficiently fast. Therefore, we have presented numerical evidence for isotropization in these models.
## III Asymptotic dynamics at early times for $`G_0`$ cosmologies
It would also be of interest to investigate general inhomogeneous ($`G_0`$) models with $`\gamma 2`$ close to the singularity. We refer to Uggla03 for the equations in the separable volume gauge, using Hubble-normalized variables (where we set $`\mathrm{\Lambda }=0`$, $`R^\alpha =0`$). We give the asymptotic decay rates below. We hope to be able to numerically simulate the $`G_0`$ cosmologies in future work.<sup>9</sup><sup>9</sup>9Simulation of $`G_0`$ cosmologies is expensive and technically difficult. Simulations of vacuum $`G_0`$ cosmologies with 50 grid points for each of $`x^i`$ has been carried out recently Garfinkle04 .
### III.1 Case $`\gamma =2`$
For the case $`\gamma =2`$, we assume that the following conditions hold uniformly for open sets of $`x^i`$.
$`C_1:`$ $`\underset{t\mathrm{}}{lim}(E_\alpha {}_{}{}^{i},A_\alpha ,N_{\alpha \beta },v_\alpha )=\mathrm{๐},\underset{t\mathrm{}}{lim}\mathrm{\Sigma }_{\alpha \beta }=\widehat{\mathrm{\Sigma }}_{\alpha \beta }`$
$`C_2:`$ $`_x(E_\alpha {}_{}{}^{i},A_\alpha ,N_{\alpha \beta },v_\alpha ,\mathrm{\Sigma }_{\alpha \beta })\text{are bounded as}t\mathrm{}.`$
$`C_3:`$ $`V=๐ช(f(t))\text{implies}_xV=๐ช(f(t))\text{(asymptotic expansions in time}`$
$`\text{can be differentiated with respect to the spatial coordinates)}.`$
Without loss of generality, we perform a spatially-dependent rotation to set $`\widehat{\mathrm{\Sigma }}_{\alpha \beta }=0`$ for $`\alpha \beta `$. The decay rates as $`t\mathrm{}`$ are then given by (no summation over repeated indices below)
$$E_\alpha {}_{}{}^{i}=e^{(2\widehat{\mathrm{\Sigma }}_{\alpha \alpha })t}[\widehat{E}_\alpha {}_{}{}^{i}+๐ช(F)],r_\alpha =e^{(2\widehat{\mathrm{\Sigma }}_{\alpha \alpha })t}[\widehat{r}_\alpha +๐ช(tF)]$$
(17)
$$A_\alpha =e^{(2\widehat{\mathrm{\Sigma }}_{\alpha \alpha })t}[\frac{1}{2}\widehat{E}_\alpha {}_{}{}^{i}_{i}^{}\widehat{\mathrm{\Sigma }}_{\alpha \alpha }t+\widehat{A}_\alpha +๐ช(tF)]$$
(18)
$$N_{\alpha \alpha }=e^{(2+2\widehat{\mathrm{\Sigma }}_{\alpha \alpha })t}[\widehat{N}_{\alpha \alpha }+๐ช(tF)]$$
(19)
$$N^{\alpha \beta }=e^{(2\widehat{\mathrm{\Sigma }}_{\mu \mu })t}[\widehat{E}_\gamma {}_{}{}^{i}_{i}^{}ฯต^{\gamma \delta (\alpha }\widehat{\mathrm{\Sigma }}{}_{}{}^{\beta )}{}_{\delta }{}^{}t+\widehat{N}^{\alpha \beta }+๐ช(tF)],\mu \alpha \beta $$
(20)
$$\mathrm{\Omega }=\widehat{\mathrm{\Omega }}+๐ช(F),v_\alpha =e^{(2\widehat{\mathrm{\Sigma }}_{\alpha \alpha })t}[\frac{1}{2}(\widehat{E}_\alpha {}_{}{}^{i}_{i}^{}\mathrm{ln}\widehat{\mathrm{\Omega }}2\widehat{r}_\alpha )t+\widehat{v}^\alpha +๐ช(tF)]$$
(21)
$$\mathrm{\Sigma }_{\alpha \alpha }=\widehat{\mathrm{\Sigma }}_{\alpha \alpha }+๐ช(F),\mathrm{\Sigma }_{\alpha \beta }=๐ช(F),\alpha \beta $$
(22)
where
$$F=t^2e^{2(2s_+)t}+e^{2(2+2s_{})t},s_+(x^i)=\underset{\alpha =1,2,3}{\mathrm{max}}\widehat{\mathrm{\Sigma }}_{\alpha \alpha },s_{}(x^i)=\underset{\alpha =1,2,3}{\mathrm{min}}\widehat{\mathrm{\Sigma }}_{\alpha \alpha }.$$
The Hubble scalar and the deceleration parameter satisfy
$$H=e^{3t}[\widehat{H}+๐ช(F)],q=2+๐ช(F).$$
(23)
The twenty eight hat variables above<sup>10</sup><sup>10</sup>10Note that $`\widehat{\mathrm{\Sigma }}_{11}+\widehat{\mathrm{\Sigma }}_{22}+\widehat{\mathrm{\Sigma }}_{33}=0`$, and we can write $`\widehat{\mathrm{\Sigma }}_{\alpha \alpha }=\text{diag}(2\widehat{\mathrm{\Sigma }}_+,\widehat{\mathrm{\Sigma }}_++\sqrt{3}\widehat{\mathrm{\Sigma }}_{},\widehat{\mathrm{\Sigma }}_+\sqrt{3}\widehat{\mathrm{\Sigma }}_{})`$. are functions of $`x^i`$, and satisfy the following sixteen constraints
$$\widehat{r}_\alpha =\widehat{E}_\alpha {}_{}{}^{i}_{i}^{}\mathrm{ln}\widehat{H},0=\widehat{E}_\alpha {}_{}{}^{i}_{i}^{}\widehat{\mathrm{\Sigma }}_{\alpha \alpha }+(2\widehat{\mathrm{\Sigma }}_{\alpha \alpha })\widehat{r}_\alpha 3\widehat{A}_\alpha \widehat{\mathrm{\Sigma }}_{\alpha \alpha }ฯต_\alpha {}_{}{}^{\beta \gamma }\widehat{N}_{\beta \delta }^{}\widehat{\mathrm{\Sigma }}_\gamma {}_{}{}^{\delta }+6\widehat{\mathrm{\Omega }}\widehat{v}_\alpha ,$$
(24)
$$\widehat{\mathrm{\Omega }}=1\frac{1}{6}(\widehat{\mathrm{\Sigma }}_{11}{}_{}{}^{2}+\widehat{\mathrm{\Sigma }}_{22}{}_{}{}^{2}+\widehat{\mathrm{\Sigma }}_{33}{}_{}{}^{2}),0=2(\widehat{E}_{[\alpha }{}_{}{}^{j}_{j}^{}\widehat{r}_{[\alpha }\widehat{A}_{[\alpha })\widehat{E}_{\beta ]}{}_{}{}^{i}ฯต_{\alpha \beta \delta }\widehat{N}^{\delta \gamma }\widehat{E}_\gamma {}_{}{}^{i}.$$
(25)
That leaves twelve hat variables, of which eight are essential, since we can use the remaining temporal gauge freedom to set $`\widehat{H}=const`$ and make a change of coordinates to set three of the $`\widehat{E}_\alpha ^i`$โs.
Note that $`\widehat{\mathrm{\Sigma }}_{\alpha \alpha }`$ satisfy $`|\widehat{\mathrm{\Sigma }}_{\alpha \alpha }|<2`$ by definition. The exponents of $`N_{\alpha \alpha }`$ must be positive, however, thus requiring that $`\widehat{\mathrm{\Sigma }}_{\alpha \alpha }>1`$; i.e., the Jacobs solutions are restricted within the triangle
$$\mathrm{\Sigma }_+<\frac{1}{2},\frac{1}{\sqrt{3}}(1+\mathrm{\Sigma }_+)<\mathrm{\Sigma }_{}<\frac{1}{\sqrt{3}}(1+\mathrm{\Sigma }_+).$$
(26)
See triangle I in Figure 3b.
The Hubble-normalized Weyl scalar $`๐ฒ`$ is given by
$$๐ฒ^2=\frac{1}{9}(\widehat{\mathrm{\Sigma }}_++\widehat{\mathrm{\Sigma }}_+^2\widehat{\mathrm{\Sigma }}_{}^2)^2+\frac{1}{9}(12\widehat{\mathrm{\Sigma }}_+)^2\widehat{\mathrm{\Sigma }}_{}^2+๐ช(tF).$$
(27)
The work of Andersson and Rendall Andersson2001 is more rigorous, but also assumes more (e.g., analyticity). Our results complement their analysis by providing more details on the spatial curvature variables and on the big O terms, and is easier to apply and interpret. We have also discussed the restrictions on the Jacobs disk.
### III.2 Case $`\gamma >2`$
For the case $`\gamma =2`$, we assume that the following conditions hold uniformly for open sets of $`x^i`$.
$`C_1:`$ $`\underset{t\mathrm{}}{lim}(E_\alpha {}_{}{}^{i},A_\alpha ,N_{\alpha \beta },v_\alpha ,\mathrm{\Sigma }_{\alpha \beta })=\mathrm{๐},`$
$`C_2:`$ $`_x(E_\alpha {}_{}{}^{i},A_\alpha ,N_{\alpha \beta },v_\alpha ,\mathrm{\Sigma }_{\alpha \beta })\text{are bounded as}t\mathrm{}.`$
$`C_3:`$ $`V=๐ช(f(t))\text{implies}_xV=๐ช(f(t))\text{(asymptotic expansions in time}`$
$`\text{can be differentiated with respect to the spatial coordinates)}.`$
The decay rates as $`t\mathrm{}`$ are then given by
$$(E_\alpha {}_{}{}^{i},r_\alpha ,A_\alpha ,N_{\alpha \beta },v_\alpha )=\eta [(\widehat{E}_\alpha {}_{}{}^{i},\widehat{r}_\alpha ,\widehat{A}_\alpha ,\widehat{N}_{\alpha \beta },\widehat{v}_\alpha )+๐ช(\xi )]$$
(28)
$$\mathrm{\Sigma }_{\alpha \beta }=\xi [\widehat{\mathrm{\Sigma }}_{\alpha \beta }+๐ช(\xi ^2)],\mathrm{\Omega }=1\frac{1}{6}\widehat{\mathrm{\Sigma }}_{\alpha \beta }\widehat{\mathrm{\Sigma }}^{\alpha \beta }\xi ^2+๐ช(\xi ^4+\eta ^2)$$
(29)
where $`\eta `$ and $`\xi `$ are given in (13). The Hubble scalar and the deceleration parameter satisfy
$$H=e^{\frac{3}{2}\gamma t}[\widehat{H}+๐ช(\xi ^2)],q=\frac{1}{2}(3\gamma 2)+๐ช(\xi ^2).$$
(30)
The hat variables above are functions of $`x^i`$, and satisfy the following constraints
$$\widehat{r}_\alpha =\widehat{E}_\alpha {}_{}{}^{i}_{i}^{}\mathrm{ln}\widehat{H},0=2\widehat{r}_\alpha +3\gamma \widehat{v}_\alpha ,0=2(\widehat{E}_{[\alpha }{}_{}{}^{j}_{j}^{}\widehat{r}_{[\alpha }\widehat{A}_{[\alpha })\widehat{E}_{\beta ]}{}_{}{}^{i}ฯต_{\alpha \beta \delta }\widehat{N}^{\delta \gamma }\widehat{E}_\gamma {}_{}{}^{i}.$$
(31)
The Hubble-normalized Weyl scalar $`๐ฒ`$ is given by
$$๐ฒ^2=\frac{1}{9}\frac{1}{6}\widehat{\mathrm{\Sigma }}_{\alpha \beta }\widehat{\mathrm{\Sigma }}^{\alpha \beta }\xi ^2+๐ช(\xi ^3).$$
(32)
We comment that the results for $`G_2`$ and $`G_0`$ cases differ slightly, due to the difference in both the temporal and spatial gauges employed.
## IV Discussion
We have shown analytically (and confirmed numerically for the $`G_2`$ models) that for the case $`\gamma =2`$, a subset of the Jacobs solutions (triangles I in Figure 3) are locally stable into the past, and for the case $`\gamma >2`$ the flat Friedmann-Lemaรฎtre solution is locally stable into the past. These results confirm and extend the BKL conjectures, and complement the work of Andersson and Rendall Andersson2001 . They are of also of current physical interest; for example, they lend support to previous isotropization results erik ; Dunsby .
In applications to specific physical theories, it may be necessary to investigate the robustness of these results in the presence of certain additional fields. For example, string and $`M`$-theories were discussed in order to motivate the existence of a scalar field (dilaton) in the early universe. However, in these theories there are also $`p`$-form fields present, which can lead to oscillatory behavior close to the cosmological singularity Damour00 .
We comment that the conditions $`C_1`$$`C_3`$ in Section II.1 can sometimes be violated in a neighbourhood of isolated timelines $`x=const`$, where formation of spiky spatial structures are observed numerically. In a preliminary numerical investigation for the stiff $`G_2`$ case, these structures appear to be similar to the spikes in vacuum models Andersson05 and in models with $`\gamma <2`$ thesis . We conjecture that the limits for $`\mathrm{\Sigma }_+`$ and $`\mathrm{\Sigma }_{}`$ along spike timelines reside in the following two triangles: (i) for true spikes, the limits reside in the triangle
$$\frac{2}{\sqrt{3}}(1+\mathrm{\Sigma }_+^H)<\mathrm{\Sigma }_{}^H<\frac{1}{\sqrt{3}}(1+\mathrm{\Sigma }_+^H),\mathrm{\Sigma }_{}^H>\sqrt{3}\mathrm{\Sigma }_+^H,$$
(33)
(ii) for false spikes, the limits reside in the triangle
$$\mathrm{\Sigma }_{}^H<\frac{1}{\sqrt{3}}(1+\mathrm{\Sigma }_+^H),\mathrm{\Sigma }_{}^H>0,\mathrm{\Sigma }_{}^H<\sqrt{3}\mathrm{\Sigma }_+^H$$
(34)
(see triangles II and III in Figure 3a respectively). We conjecture that true spikes are physically real, and are reflected by a discontinuous limit for the Weyl scalar. On the other hand, false spikes are artifacts of the rotating frame, and are reflected by a continuous limit for the Weyl scalar. It is of interest to study these spikes in more detail.
Similarly, for the stiff $`G_0`$ case we conjecture that spikes occur when the limits of $`\mathrm{\Sigma }_{\alpha \alpha }`$ are discontinuous. When this occurs, the limits are confined to the three triangles
$$0<4+2\mathrm{\Sigma }_{\alpha \alpha }\mathrm{\Sigma }_{\beta \beta },\widehat{\mathrm{\Sigma }}_{\alpha \alpha }<1,$$
(35)
for all $`\alpha \beta `$ (see triangles II in Figure 3b; also compare with (WainwrightHsu89, , Figure 3)).
###### Acknowledgements.
This work was supported by NSERC of Canada. We would like to thank Henk van Elst, Claes Uggla and Alan Rendall for their helpful comments.
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# Experimental motivation and empirical consistency in minimal no-collapse quantum mechanics
## I Introduction
Historically, quantum theory was motivated by the need to describe the behavior of microscopic systems not explainable by the laws of classical physics. Not only was quantum mechanics deemed unnecessary for a description of the macroworld of our experience, it also led to โstrangeโ consequences that seemed to blatantly contradict our experience, as famously illustrated by the Schrรถdinger-cat Gedanken experiment Schrรถdinger (1935) and later generally referred to as the โmeasurement problem.โ Therefore quantum theory was often, as in the Copenhagen interpretation, banned a priori from the macrosopic realm.
Over the past decade, however, a rapidly growing number of experiments have demonstrated the existence of quantum superpositions of mesoscopically and macrosopically distinct states on increasingly large scales. Such superpositions are observed as individual quantum states and are perfectly explained by unitarily evolving wave functions. On the other hand, decoherence theory Zeh (1970, 1973); Zurek (1981, 2003a); Joos et al. (2003); Schlosshauer (2004) has enabled one to understand the fragility of such superpositions, and thus the extreme difficulty in observing them outside of sophisticated experimental setups, as being due to ubiquitous quantum interactions with environmental degrees of freedom.
These developments have thus extended the domain for an application of quantum theory far into the mesoscopic and macroscopic realm, which lends strong support to assuming a universally exact and applicable Schrรถdinger equation. To make a physically compelling case for such a purely unitary quantum theory we must pursue two related goals. First, we ought to continue to design experiments which demonstrate the existence of quantum superpositions of macrosopically distinct states โ and which, ideally, can explicitly rule out collapse models. Second, since the assumption of a universal Schrรถdinger dynamics implies that superpositions of (presumably macroscopically) different observer states are both possible and inescapable if we include physical observers into the quantum-mechanical description, we must simultaneously show that environmental decoherence provides the necessary and sufficient mechanism to explain our observation of a โclassicalโ world. The emergence of the latter can then be understood not only in spite of, but precisely because of the quantum formalism โ no classical prejudice need to be imposed.
The formal basis for a derivation of a viable interpretation of quantum mechanics from the โbareโ unitary formalism alone has been outlined in several papers. The basic idea was introduced in Everettโs proposal of a relative-state view of quantum mechanics Everett (1957). It was later adapted and popularized by deWitt DeWitt (1970, 1971); DeWitt and Graham (1973) in his โmany-worldsโ interpretation of quantum mechanics, whose elements go far beyond the abstract sketches of Everett und which must therefore be strictly distinguished from Everettโs proposal Kent (1990). Relative-state interpretations were subsequently fleshed out, by taking into account decoherence effects, in works by Zeh Zeh (1970, 1973, 2000), Zurek Zurek (1998, 2003a, 2004a), Wallace Wallace (2003a, 2002), and others (see, for example, Deutsch (1985); Vaidman (1998); Donald (1999)). Such a theory can be based on the most minimal set of assumptions about the quantum formalism and its interpretation. First, a completely known (pure) state of an isolated quantum system $`๐ฎ`$ is described by a normalized state vector $`|\psi `$ in a Hilbert space $`H_๐ฎ`$. Second. the time evolution of a state vector $`|\psi `$ is given by the Schrรถdinger equation $`\mathrm{i}\mathrm{}\frac{}{t}|\psi =\widehat{H}_๐ฎ`$, where $`\widehat{H}_๐ฎ`$ is the Hamiltonian of the system $`๐ฎ`$. No mention is made of measurements in this formulation. Instead, measurements are described without special axioms in terms of physical interactions between systems described by state vectors (wave functions) and governed by suitable interaction Hamiltonians. Observables then emerge as a derived concept (see, for example, Joos et al. (2003); Zurek (2003a)).
In this paper, however, we take a less formal route and focus on an analysis of the experimental and theoretical progress (with an emphasis on the former) towards the two goals mentioned before, namely, the continued acquisition of experimental evidence for superpositions of macrosopically distinct states and an explanation for the emergence of definite perceptions in spite of an assumed universal validity of the superposition principle.
Our goal is to show that there is no experimental evidence for a breakdown of the superposition principle and the related interference effects at any length scale investigated thus far. Whenever a decay of such superpositions is observed, it can be fully accounted for (both experimentally and theoretically) as resulting from environmental interactions. The absense of any empirical evidence for nonlinear deviations from unitary time evolution, combined with the ability to give an empirically adequate description of the decoherence of superpositions into apparent mixtures, provides good reasons to take the universal validity of the Schrรถdinger equation as a working assumption and to explore the consequences of this assumption.
The resulting theory will require more attention to a detailed quantum-mechanical description of observers and observations. Such an account is interpretation-neutral, while the question of its relevance for solving the measurement problem may depend on the particular features of an interpretation. This is so because there exist interpretations, for example, Bohmian mechanics or modal interpretations, that claim to solve the measurement problem *without* having to give an explicit account of the physical processes describing observers and observations (see also Sec. III).
This paper is organized as follows. In Sec. II, we shall discuss and analyze three important experimental domainsโsuperconducting quantum interference devices (SQUIDs), matter-wave interferometry, and Bose-Einstein condensationโthat have provided evidence for superpositions of macroscopically distinguishable states. Sec. III comments on the current status of physical collapse theories in view of the described experiments. In Secs. IV and V, we shall discuss steps towards the resolution of two issues that have often been considered as posing a challenge to relative-state interpretations: The question of the origin of quantum probabilities and the connection with Bornโs rule, and the problem of the โobjectificationโ of observables and thus the emergence of โclassical reality.โ Sec. VI analyzes theoretical models for decoherence in the perceptive and cognitive apparatus, and the implications of such decoherence processes. Finally, in Sec. VII, we shall summarize our main conclusions and discuss possible next steps.
## II Superpositions of macroscopically distinct states: Experiments and implications
In the following, we shall describe three recent experimental areas that have led to (or that are very close to achieving) the observation of superpositions of mesoscopically and macroscopically distinguishable states: Coherent quantum tunneling in SQUIDs (Sec. II.2), diffraction of C<sub>70</sub> (and larger) molecules in matter-wave interferometers (Sec. II.3), and number-difference superpositions in two-species Bose-Einstein condensates (Sec. II.4). These experiments have achieved the largest such superpositions observed thus far and also represent the most promising experimental domains for achieving even larger superpositions in the future.
For some earlier experiments demonstrating mesoscopic and macrosopic quantum effects, see the setups using superconductors Clarke et al. (1988); Rouse et al. (1995); Silvestrini et al. (1997); Rouse et al. (1998); Nakamura et al. (1999), nanoscale magnets Friedman et al. (1996); Wernsdorfer et al. (1997); Barco et al. (1999), laser-cooled trapped ions Monroe et al. (1996), and photons in a microwave cavity Brune et al. (1996); Raimond et al. (1997). We would also like to mention Leggettโs review article Leggett (2002) which discusses some experiments that probe the limits of quantum mechanics. Leggettโs motivation, however, is somewhat different than that of the present author, as Leggettโs main aim is to assess the status of physical collapse theories in view of these experiments.
### II.1 Measuring the macrosopic distinctness of states in a superposition
Before embarking on an analysis of the experiments, we shall first lend a more precise meaning to the ubiquitous phrase โsuperposition of macrosopically distinct (or distinguishable) states.โ If confronted with a superposition of two states $`|A`$ and $`|B`$ of the form
$$|\mathrm{\Psi }=\frac{1}{\sqrt{2}}\left(|A+|B\right),$$
(1)
how are we to decide whether this indeed represents a macrosopic Schrรถdinger-cat state? Clearly, two conditions will need to be fulfilled:
1. The states $`|A`$ and $`|B`$ must differ macrosopically in some extensive quantity (e.g., spatial separation, total mass, magnetic moment, momentum, charge, current, etc.), relative to a suitable microsopic reference value.
2. The degree of GHZ-type entanglement Greenberger et al. (1990) in the state $`|\mathrm{\Psi }`$, i.e., the number of correlations that would need to be measured in order to distinguish this state from a mixture, must be sufficiently large. With $`|A`$ and $`|B`$ usually representing GHZ-like multi-particle states in complex systems such as superconducting currents, molecules, and atomic gases, this measure can typically be well-estimated by the number of microsopic constituents (electrons, protons, neutrons) in the system.
A similiar combination of two measures has been suggested by Leggett Leggett (1980, 2002) under the labels โextensive differenceโ and โdisconnectivity.โ We shall adopt Leggettโs former term for the first condition, and use the term โdegree of entanglementโ for the second. A both necessary and sufficient condition for a superposition to be considered a superposition of macroscopically distinct states is then given by the requirement that both the extensive difference between $`|A`$ and $`|B`$ and the interparticle entanglement in $`|\mathrm{\Psi }`$ be large relative to an appropriate microsopic unit.
### II.2 Superconducting quantum interference devices
Experiments using SQUIDs have not only demonstrated that the dynamics of a macrosopic quantity of matter (here $`10^9`$ Cooper pairs) can be collectively determined by a single macrosopic coordinate governed by quantum mechanics, but have also achieved the creation and indirect observation of quantum superpositions of two truly macrosopic states that correspond to currents of several $`\mu `$A running in opposite directions.
#### II.2.1 SQUID setup and detection of superpositions of macroscopically distinct currents
A SQUID consists of a superconducting loop interrupted by a Josephson junction and immersed into an external magnetic field that creates a flux $`\mathrm{\Phi }_{\text{ext}}`$ through the loop. This allows for a persistent dissipationless current (โsupercurrentโ) to flow around the loop, in clockwise or counterclockwise direction, creating an additional flux. Such a current is composed of a very large number of Cooper pairs (i.e., Bose-condensed electron pairs) whose collective center-of-mass motion can be described by a macrosopic wave function around the loop.
Since the wave function must be continuous around the loop, an integer $`k`$ times its wavelength must equal the circumference of the loop. Since the Josephson junction induces a discontinuous phase drop $`\mathrm{\Delta }\varphi _J`$, and since the total change in phase around the superconducting loop is given by $`2\pi \mathrm{\Phi }/\mathrm{\Phi }_0`$, where $`\mathrm{\Phi }_0=h/2e`$ is the flux quantum and $`\mathrm{\Phi }`$ is the total trapped flux through the loop, the phase continuity condition implies
$$\mathrm{\Delta }\varphi _J+2\pi \mathrm{\Phi }/\mathrm{\Phi }_0=2\pi k,$$
(2)
with $`k=1,2,\mathrm{}`$. This means that the collective quantum dynamics of the SQUID are determined by the single macrosopic variable $`\mathrm{\Phi }`$.
The effective SQUID Hamiltonian can be written as Weiss (1999)
$$\begin{array}{c}\widehat{H}=\frac{\widehat{P}_\mathrm{\Phi }^2}{2C}+U(\mathrm{\Phi })=\frac{\mathrm{}^2}{2C}\frac{d^2}{d\mathrm{\Phi }^2}+[\frac{(\mathrm{\Phi }\mathrm{\Phi }_{\text{ext}})^2}{2L}\hfill \\ \hfill \frac{I_c\mathrm{\Phi }_0}{2\pi }\mathrm{cos}\left(2\pi \frac{\mathrm{\Phi }}{\mathrm{\Phi }_0}\right)],\end{array}$$
(3)
where $`C`$ is the total capacitance (mainly due to the junction), $`L`$ is the (finite) self-inductance of the loop, and $`I_c`$ is the critical current of the junction. This Hamiltonian induces dynamics that are analogous to the motion of a particle with effective โmassโ $`C`$ moving in $`\mathrm{\Phi }`$-space in a tilted one-dimensional double-well potential, with the tilt determined by $`\mathrm{\Phi }_{\text{ext}}`$. The role of the canonical variables $`\widehat{X}`$ and $`\widehat{P}`$ is here played by the total trapped flux $`\widehat{\mathrm{\Phi }}`$ and the total displacement current $`\widehat{P}_\mathrm{\Phi }=\mathrm{i}\mathrm{}d/d\widehat{\mathrm{\Phi }}`$ (which has units of charge; $`Cd\widehat{P}_\mathrm{\Phi }/dt`$ is the charge difference across the junction).
A set of eigenstates $`|k`$ of the Hamiltonian of Eq. (3), called โ$`k`$-fluxoid states,โ are localized in one of the wells of the potential below the (classically impenetrable) barrier if the damping induced by the Josephson junction is weak. The corresponding wave functions $`\psi _k(\mathrm{\Phi })\mathrm{\Phi }|k`$ are locally $`s`$-harmonic, so their amplitudes are peaked around the respective minimum of $`U(\mathrm{\Phi })`$ with narrow spreads in flux space. Thus these low-lying energy eigenstates have a relatively small range of associated flux values and can therefore (at least for sufficiently small $`k`$) also be viewed as โfuzzyโ eigenstates of the flux operator. By adjusting $`\mathrm{\Phi }_{\text{ext}}`$, the energy levels are shifted, and for certain values of $`\mathrm{\Phi }_{\text{ext}}`$, two levels in opposite wells can be made to align, which allows for resonant quantum tunneling between the wells (i.e., between two fluxoid states) Silvestrini et al. (1996); Rouse et al. (1998), leading to a macroscopic change in the magnetic moment of the system.
The most important states for our subsequent treatment are the zero-fluxoid state $`|0`$ and the one-fluxoid state $`|1`$. Since the states $`|0`$ and $`|1`$ are localized in, respectively, the left and right well of the potential, let us denote them by $`|L`$ and $`|R`$ in the following. These states correspond (apart from the quantum zero-point energy van der Wal et al. (2000)) to a classical persistent-current state and thus to macrosopically distinguishable directions of the superconducting current. Since other states are well-separated in energy, the SQUID can thus be effectively modelled as a macroscopic quantum-mechanical two-state system (i.e., as a macrosopic qubit).
At bias $`\mathrm{\Phi }_{\text{ext}}=\mathrm{\Phi }_0/2`$, the well becomes symmetric and the corresponding two fluxoid states $`|L`$ and $`|R`$ would become degenerate (see Fig. 1). However, the degeneracy is lifted by the formation of symmetric and antisymmetric superpositions of $`|L`$ and $`|R`$ that represent the new energy ground state,
$$|\mathrm{\Psi }_s=\frac{1}{\sqrt{2}}\left(|L+|R\right)$$
(4)
with energy $`E_+`$, and the first excited energy eigenstate
$$|\mathrm{\Psi }_a=\frac{1}{\sqrt{2}}\left(|L|R\right)$$
(5)
with energy $`E_{}`$. Thus these eigenstates are delocalized across the two wells. The (typically very small) energy splitting $`\mathrm{\Delta }E=E_aE_s`$ is determined by the WKB matrix elements for tunneling between the two wells (and thus between $`|L`$ and $`|R`$), and is only dependent on the capacitance $`C`$ of the junction, scaling as $`\mathrm{\Delta }E\mathrm{e}^\sqrt{C}`$.
If the system is now more generally described by an arbitrary superposition of $`|L`$ and $`|R`$, $`|\mathrm{\Psi }(t)=\alpha (t)|L+\beta (t)|R`$, and if we choose the left-localized state $`|L`$ as the initial state of the SQUID, i.e., $`|\mathrm{\Psi }(t=0)=|L`$, we obtain the time evolution
$$|\mathrm{\Psi }(t)|L\mathrm{cos}(\mathrm{\Delta }Et/2)+i|R\mathrm{sin}(\mathrm{\Delta }Et/2).$$
(6)
Thus the wave function oscillates coherently between the two localized current states $`|L`$ and $`|R`$ in each well (see Fig. 2) at a rate determined by $`\mathrm{\Delta }E`$, since the probability to find the wave function localized in, say, the left well is oscillatory in time,
$$P_L(t)=|L|\mathrm{\Psi }(t)|^2=\mathrm{cos}^2(\mathrm{\Delta }Et/2).$$
(7)
This leads to coherent quantum tunneling between the two wells and manifests itself in an oscillation of the current in the SQUID between clockwise and counterclockwise directions. This tunneling effect has been directly observed in superconducting qubit setups similiar to the one described here Nakamura et al. (1999); Korotkov and Averin (2001); Korotkov (2001); Greenberg et al. (2002); Martinis et al. (2002); Yu et al. (2002); Vion et al. (2002).
The indirect route for detecting the presence of superpositions of states corresponding to macrosopic currents running in opposite directions relies on a static spectroscopic measurement of the energy difference $`\mathrm{\Delta }E`$ (see Fig. 1). Friedman *et al.* Friedman et al. (2000) have confirmed the existence of such an energy gap (in excellent agreement with theoretical predictions) and, therefore, of superpositions of macroscopically distinct fluxoid states (see also van der Wal et al. (2000) for a similiar experiment and result). In their setup, $`|L`$ and $`|R`$ (which in this experiment corresponded to $`k=4`$ and $`k=10`$, respectively) differed in flux by more than $`\mathrm{\Phi }_0/4`$ and in current by 2โ3 $`\mu `$A, corresponding to about $`10^{10}\mu _B`$ in local magnetic moment. Furthermore, the dynamics of the in-unison motion of the approximately $`10^9`$ Cooper pairs represented by $`|L`$ and $`|R`$ are given by a single unitarily evolving wave function representing the collective flux coordinate $`\mathrm{\Phi }`$.
#### II.2.2 Scaling
A main advantage of SQUIDs over other experiments (such as those described in the subsequent sections) that probe the limits of quantum mechanics lies in the fact that the relevant macrosopic variable, namely, the trapped flux through the SQUID ring, can be controlled by means of microsopic energy differences in the Josephson junction Leggett (2002). As mentioned before, the tunneling matrix element scales as $`\mathrm{e}^\sqrt{C}`$, where $`C`$ is dominantly determined by the junction rather than by the size of the loop. Thus the difficulty of observing superpositions of macrosopically distinct states scales essentially independently of the degree of macrosopic distinctness between these states (i.e., difference in flux between the opposite currents). This is in stark contrast to the matter-wave diffraction experiments and Bose-Einstein condensates discussed below. In the first case, the grating spacing must decrease as $`1/\sqrt{N}`$ with the number $`N`$ of atoms in the molecule, in the second case the decoherence rate increases as $`N^2`$ with the number $`N`$ of atoms in the condensate.
This particular property of SQUIDs has allowed for the creation of superpositions of states that differ by several orders of magnitude more than in other experiments (see Sec. II.5 below).
#### II.2.3 The interpretation of superpositions
It is well known that quantum-mechanical superpositions must not be interpreted as a simple superposition (addition) of probability distributions. Formally, this conclusion is of course well-reflected in the fact that, in quantum mechanics, we deal with superpositions of probability amplitudes rather than of probabilities, leading to interference terms in the probability distribution.
However, this crucial difference between classical and quantum-mechanical superpositions is sometimes not sufficiently clearly brought out when describing particular experimental situations. In the case of the standard double-slit interference experiment, for example, the state of the diffracted particle is described by a coherent superposition $`|\psi =\left(|\psi _1+|\psi _2\right)/\sqrt{2}`$ of the states $`|\psi _1`$ and $`|\psi _2`$ corresponding to passage through slit 1 and 2, respectively. This is frequently interpreted as simply representing simultaneous passage of the particle through both slits, i.e., presence of the particle in two distinct spatial regions at the same time, thereby tacitly neglecting the interference terms in the probability distribution.
In the double-slit example, this view will not necessarily be disproven until the stage of the screen is reached at which interference fringes appear. Similiarly, and even more drastically, the superpositions of macrosopically distinct current states in a SQUID show that the simplified view of a classical superposition of probability distributions is inadequate. For, if this view were correct, the two contributing opposite currents would mutually cancel out and thus the net โcurrentโ described by this superposition would have to be zero, contrary to what is observed. Instead, the SQUID opposite-current superposition represents a novel individually existing physical state that can be described as a coherent โinteractionโ between simultaneously present states representing currents of opposite direction.
The SQUID example also shows that the โsplittingโ often referred to in an Everettian framework (for example, in deWittโs popularization of the โmany-worlds viewโ DeWitt (1970, 1971); DeWitt and Graham (1973)) should not be taken too literally. The transition, i.e., the โsplit,โ from a single โclassicalโ stateโi.e., classically defined definite structures such as particles (defined as having a definite position), currents (defined as a flow of charge into a definite direction), etc.โinto a state describing a superposition of such states occurs in a completely unitary and thus reversible manner by changing $`\mathrm{\Phi }_{\text{ext}}`$. There is only one single global state vector $`|\mathrm{\Psi }(t)`$ at all times that corresponds to โphysical reality.โ The decomposition into a superposition of other states is a primarily formal procedure useful in revealing the physical quantities of our experience contained in the arbitrary state vector $`|\mathrm{\Psi }(t)`$, since the latter can in general not be related to any โclassicalโ physical structure that would correspond to directly observed objects or properties. In this sense, the โsplitโ is simply a consequence of trying to trace throughout time a particular (usually โclassicalโ) state that does not coincide with $`|\mathrm{\Psi }(t)`$. Quantum mechanics shows that this can, in general, only be done in a relative-state sense.
The decomposition obtains also physical meaning when the dynamical evolution of the system described by $`|\mathrm{\Psi }(t)`$ is considered, as the coefficients multiplying the โclassicalโ terms in the superposition will in general be time-dependent. In the example of the SQUID, the coherent-tunneling state does not directly relate to a current in the classical sense (i.e., a current of definite direction), but it can be decomposed into two such currents of opposite direction. The physical relevance of this decomposition and the meaning of the superposition then manifests itself as a current that oscillates between clockwise and counterclockwise directions.
#### II.2.4 Decoherence and the preferred basis
A particularly interesting feature of the macrocurrent superpositions in SQUIDs is the fact that the interaction with the environment leads to a localization in flux space, rather than to the much more familiar and common localization in position space. In other words, the โpreferred basisโ (Zurekโs โpointer statesโ Zurek (1981, 1982)) of the SQUID are flux eigenstates.
This observation is perfectly well accounted for by decoherence theory, which describes the selection of the preferred basis by means of the stability criterion, first formulated by Zurek Zurek (1981) (see also Zurek (1982, 1993, 1998, 2003a); Schlosshauer (2004)). According to this criterion, the basis used to represent the possible states of the system must allow for the formation of dynamically stable system-environment correlations. A sufficient (albeit not necessary) requirement for this criterion to be fulfilled is given by the condition that all basis projectors $`\widehat{P}_n=|s_ns_n|`$ of the system must (at least approximately) commute with the system-environment interaction Hamiltonian $`\widehat{H}_{\text{int}}`$, i.e.,
$$[\widehat{H}_{\text{int}},\widehat{P}_n]=0\text{for all }n\text{.}$$
(8)
That is, the preferred basis of the system is given by a set of eigenvectors of $`\widehat{H}_{\text{int}}`$.
In the case of the SQUID experiments at bias $`\mathrm{\Phi }_{\text{ext}}=\mathrm{\Phi }_0/2`$, if the interaction with the environment is very weak and thus the dynamics of the SQUID system are dominantly governed by the effective SQUID Hamiltonian $`\widehat{H}`$, Eq. (3), the preferred states are predicted to be eigenstates of this Hamiltonian, namely, the dislocalized coherent superpositions $`|\mathrm{\Psi }_s=\frac{1}{\sqrt{2}}\left(|L+|R\right)`$ and $`|\mathrm{\Psi }_a=\frac{1}{\sqrt{2}}\left(|L|R\right)`$ of the localized zero-fluxoid and one-fluxoid states $`|L`$ and $`|R`$. This is in agreement both with the observation of coherent quantum tunneling between the wells and with the measurement of the energy gap $`\mathrm{\Delta }E=E_aE_s`$ between the states $`|\mathrm{\Psi }_s`$ and $`|\mathrm{\Psi }_a`$.
Under realistic circumstances, however, the SQUID is coupled to a dissipative environment $``$ which can quite generally be modeled as a harmonic heat bath of bosons Weiss (1999), i.e., as a bath of $`N`$ harmonic oscillators with generalized coordinates $`x_\alpha `$ and $`p_\alpha `$, natural frequency $`\omega _\alpha `$, mass $`m_\alpha `$, and Hamiltonian
$$\widehat{H}_{}=\frac{1}{2}\underset{\alpha =1}{\overset{N}{}}\left(\frac{p_\alpha ^2}{m_\alpha }+m_\alpha \omega _\alpha ^2x_\alpha ^2\right).$$
(9)
The reservoir modes $`x_\alpha `$ couple dynamically to the total flux variable $`\mathrm{\Phi }`$ of the SQUID ring. More precisely, they couple to the fluxoid (and essentially opposite-current) states $`|L`$ and $`|R`$ via the interaction Hamiltonian Weiss (1999)
$$\widehat{H}_{\text{int}}=\sigma _z\left(\frac{\phi _0}{2}\underset{\alpha }{}c_\alpha x_\alpha \right),$$
(10)
where $`\sigma _z=\left(|LL||RR|\right)`$ is the so-called โpseudospinโ operator (owing its name to the fact that the SQUID double-well system can be effectively mapped onto a two-state spin system, with $`|L`$ and $`|R`$ corresponding to, say, spin โupโ and โdown,โ respectively), and $`\pm \phi _0`$ are the flux values associated with the two localized states $`|L`$ and $`|R`$.
According to the commutativity criterion, Eq. (8), the stable states into which the system decoheres are then eigenstates of $`\sigma _z`$, i.e., the preferred basis of the system is given by the two states $`|L`$ and $`|R`$, This, of course, is in full agreement with observations and explains the localization in flux space, i.e., the rapid reduction of the superposition into an apparent ensemble of the macroscopically distinguishable current states $`|L`$ and $`|R`$.
Fig. 3 illustrates this gradual disappearance of interference in the symmetric ground state $`|\mathrm{\Psi }_s=\frac{1}{\sqrt{2}}\left(|L+|R\right)`$ due to the interaction of the SQUID ring with a dissipative thermal bath in the Wigner representation of the local density operator of the SQUID Everitt et al. (2004) (see also Chudnovsky and Kuklov (2003)). As predicted by the stability criterion, the robust states (i.e., the preferred basis) selected by the environment are the macroscopically distinguishable current states $`|L`$ and $`|R`$. The resulting local loss of coherenceโthat is, the distribution of coherence, initially associated with the SQUID, over the many degrees of freedom of the SQUID-environment combinationโconstitutes the main obstacle in the observation of coherent quantum tunneling.
### II.3 Molecular matter-wave interferometry
Recent experiments by the group of Zeilinger *et al.* Arndt et al. (1999); Brezger et al. (2002); Hackermรผller et al. (2003a); Arndt et al. (2002); Nairz et al. (2003); Hornberger et al. (2003); Hackermรผller et al. (2003b, 2004); Hornberger et al. (2005) have pushed the boundary for the observation of quantum (โwaveโ) behavior towards larger and larger particles. In the experiment to be described, mesoscopic C<sub>60</sub> molecules (so-called fullerenes) and C<sub>70</sub> molecules have been observed to lead to an interference pattern following passage through a diffraction grating (โmatter-wave interferometryโ). The carbon atoms in the C<sub>70</sub> molecule are arranged in the shape of an elongated buckyball with a diameter of about 1 nm (see Fig. 4). They are complex and massive enough to exhibit properties that position them in the realm of classical solid objects rather than that of atoms. For example, they possess a large number of highly excited internal rotational and vibrational degrees of freedom that allow one to attribute a finite temperature to each individual molecule, and heated C<sub>70</sub> molecules are observed to emit photons and electrons. The particle aspect seems to be overwhelmingly clear, and yet these molecules have been shown to exhibit quantum interference effects.
#### II.3.1 Experimental setup and observation of interference
The observation of C<sub>70</sub> interference patterns and their controlled disappareance due to environmental decoherence induced by various sources has been made possible by the so-called Talbot-Lau interferometer Brezger et al. (2002) that has two main advantages over earlier setups used for molecular interferometry Bordรฉ et al. (1994); Chapman et al. (1995). First, the incident beam of molecules does not need to be collimated, allowing for much higher transmitted intensities. Second, the required period of the gratings used to obtain the interference pattern scales only with the square root of the de Broglie wavelength of the molecules, allowing for the probing of the quantum behavior of, say, sixteen times larger molecules by using an only four times smaller grating spacing.
The Talbot-Lau effect is based on the fact that the transverse part of a plane wave $`\psi (z)=\mathrm{e}^{ikz}`$ incident on a periodic grating located in the $`xy`$ plane will be identical to the grating pattern at integer multiples of the distance (โTalbot lengthโ)
$$L_\lambda =\frac{d^2}{\lambda }$$
(11)
behind the grating. Since this is a pure interference effect, the presence of the grating pattern at multiples of the Talbot length indicates the wave nature of the incident beam.
The experimental setup that makes use of the Talbot-Lau effect is shown schematically in Fig. 5. The main part consists of a set of three gold gratings with a period of about $`d=1`$ $`\mu `$m. The first grating acts as a collimator that induces a sufficient degree of coherence in the incident uncollimated beam of C<sub>70</sub> molecules in order to approximate the plane-wave assumption made above. Each point of the grating can then be viewed as representing a narrow source. The velocity of the molecules can be selected over a range from about 80 m/s to 220 m/s, corresponding to de Broglie wavelengths of approximately 2โ6 pm. The second grating is the actual diffraction element, assuming the role of the single grating in the above plane-wave example. The third grating, placed behind the second grating at a distance $`L`$ equal to the Talbot length $`L_{\lambda _{\text{C}70}}=d^2/\lambda _{\text{C}70}`$, where $`\lambda _{\text{C}70}`$ is the de Broglie wavelength of the molecules, can be moved in the $`x`$-direction and serves as a scanning detection mask for the molecular density pattern in the transverse plane at this location. The molecules that have passed through the third grating are ionized by a laser beam and then counted by an ion detector.
If the C<sub>70</sub> molecules indeed possess a quantum-wave nature, the Talbot-Lau effect implies that the molecular density pattern at the position of the third grating should consist of interference fringes with a period equal to the spacing $`d`$ of the grating pattern. Thus, when the third grating is scanned in the $`x`$-direction, we expect an oscillation in the number of transmitted molecules with period $`d`$. This is indeed what has been observed experimentally Arndt et al. (1999, 2002); Brezger et al. (2002); Nairz et al. (2003); Hornberger et al. (2003) (Fig. 6). The possibility that these fringes could result from a classical blocking of rays by the gratings (Moirรฉ fringes) can be excluded, because such patterns would be independent of the de Broglie wavelength, in contrast to what is observed experimentally Brezger et al. (2002); Hackermรผller et al. (2003b). This confirms the quantum origin of the measured fringes and thus the wave nature of the C<sub>70</sub> molecules.
It should be emphasized that the fringes represent single-particle interference effects, rather than being due to interference between different molecules Nairz et al. (2003). The latter case would require the interfering molecules to be in the same state, which is practically never the case due to the large number of different excited internal states. Furthermore, the density in the molecular beam is relatively low, such that the average distance between two molecules is much greater than the range of any intermolecular force. Thus, even if the molecules passed at such a slow rate through the apparatus that only a single molecule was present at any time, an interference pattern would emerge. The interference effect is entirely due to the splitting and overlapping of the wave fronts associated with each individual C<sub>70</sub> molecule. This demonstrates clearly that quantum-mechanical superpositions in configuration space describe individual states that can exhibit interference effects (i.e., phase dependencies) without any statistical aspect.
#### II.3.2 Disappearance of interference due to controlled decoherence
General numerical estimates for decoherence rates derived from theoretical expressions Joos and Zeh (1985); Gallis and Fleming (1990); Tegmark (1993); Hornberger and Sipe (2003) have clearly demonstrated the extreme efficiency of decoherence on mesoscopic and macrosopic scales. It is therefore usually practically impossible to control the environment in such a way as to explicitely resolve and observe the gradual action of decoherence on larger objects.
The Talbot-Lau interferometer, however, has made such observations possible and has also led to direct confirmations of the predictions of decoherence theory for mesoscopic objects Hackermรผller et al. (2003b); Hornberger et al. (2003); Hackermรผller et al. (2004); Hornberger et al. (2004, 2005). The main sources of decoherence that have been experimentally investigated are collisions with gas molecules present in the interferometer Hornberger et al. (2003); Hackermรผller et al. (2003b); Hornberger et al. (2004), and thermal emission of radiation when the C<sub>70</sub> molecules are heated to temperatures beyond 1,000 K Hackermรผller et al. (2004); Hornberger et al. (2005). Here, we shall focus on the first case of decoherence, as collisions with environmental particles represent the most natural and ubiquituous source of decoherence in nature.
In the experiments, the vacuum chamber containing the interferometer is filled with gases at different pressures. Each collision between a gas particle and a C<sub>70</sub> molecule entangles their motional states. Since the C<sub>70</sub> molecules are much more massive than the gas molecules, the motional state of the gas molecule is distinguishably changed in the collision, while the motion of the C<sub>70</sub> molecule remains essentially unaffected and can therefore still be detected at the third grating. Thus, each collision encodes which-path information about the trajectory of the C<sub>70</sub> molecule in the environment (i.e., in the colliding gas particle). This leads to decoherence in the spatial wave function of the C<sub>70</sub> molecules, since the post-collision environmental states are approximately orthogonal in the position basis due to the significant change of the motional state of the gas molecules in the collisions.
To see this more explicitely, let us denote the state of the C<sub>70</sub> molecule before and after the scattering by
$`|\psi _{\text{C}70}`$ $`={\displaystyle ๐๐ฑ\left(๐ฑ|\psi \right)_{\text{C}70}|๐ฑ_{\text{C}70}}`$ (12)
and
$`|\psi ^{}_{\text{C}70}`$ $`={\displaystyle ๐๐ฑ\left(๐ฑ|\psi ^{}\right)_{\text{C}70}|๐ฑ_{\text{C}70}},`$ (13)
respectively, where
$$\left(๐ฑ|\psi \right)_{\text{C}70}\left(๐ฑ|\psi ^{}\right)_{\text{C}70}$$
(14)
for all $`๐ฑ`$. A collision at $`๐`$ changes the state of the colliding gas molecule from $`|\phi _{\text{gas}}`$ to $`|\phi ^{},๐_{\text{gas}}`$, which encodes which-path information about the C<sub>70</sub> molecule. Since the $`|\phi ^{},๐_{\text{gas}}`$ represent distinguishable motional states, the environmental states corresponding to scattering events at different locations become approximately orthogonal,
$$\left(\phi ^{},๐|\phi ^{},๐\right)_{\text{gas}}\delta (๐๐).$$
(15)
The collision leads to an entangled state for the combined gas-C<sub>70</sub> system,
$$\begin{array}{c}|\mathrm{\Psi }_0=|\psi _{\text{C}70}|\phi _{\text{gas}}\hfill \\ \hfill |\mathrm{\Psi }๐๐\left(๐|\psi \right)_{\text{C}70}|๐_{\text{C}70}|\phi ^{},๐_{\text{gas}}.\end{array}$$
(16)
The reduced density matrix for the C<sub>70</sub> molecule expressed in the position basis is then obtained by averaging over all possible states $`|\phi ^{},๐_{\text{gas}}`$ of the gas molecule,
$`\rho _{\text{C}70}`$ $``$ $`{\displaystyle ๐๐๐๐^{}๐๐^{\prime \prime }\left(๐|\psi \right)_{\text{C}70}\left(๐^{}|\psi \right)_{\text{C}70}^{}}`$ (17)
$`\times \left(\phi ^{},๐^{\prime \prime }|\phi ^{},๐\right)_{\text{gas}}`$
$`\times \left(\phi ^{},๐^{}|\phi ^{},๐^{\prime \prime }\right)_{\text{gas}}(|๐๐^{}|)_{\text{C}70}`$
$``$ $`{\displaystyle ๐๐\left|\left(๐|\psi \right)_{\text{C}70}\right|^2\left(|๐๐|\right)_{\text{C}70}},`$
where the vanishing of interference terms $`\left(๐|\psi \right)_{\text{C}70}\left(๐^{}|\psi \right)_{\text{C}70}^{}`$, $`๐๐^{}`$, in the last step follows from the approximate orthogonality of the $`|\phi ^{},๐_{\text{gas}}`$. Thus, the gas molecules carry away which-path information, leading to a diffusion of coherence into the environment. Incidentally, in this sense, Bohrโs complementarity principle can be understood as a consequence of entanglement: The observability of an interference pattern, and thus the degree of the โwave aspectโ of the C<sub>70</sub> molecules, is directly related to the amount of information, encoded through entanglement with the state of the gas particles, about the path (the โparticle aspectโ) of the molecules.
We expect the visibility $`V_\lambda `$ of the interference fringes (defined as $`(c_{\text{max}}c_{\text{min}})/(c_{\text{max}}+c_{\text{min}})`$, where $`c_{\text{max}}`$ and $`c_{\text{min}}`$ are the maximum and minimum amplitudes of the interference pattern) to decrease as the pressure of the environmental gas is increased. A theoretical analysis Hackermรผller et al. (2003b); Hornberger and Sipe (2003); Hornberger et al. (2004) predicts that $`V_\lambda `$ will decrease exponentially with the pressure $`p=nk_BT`$ of the colliding gas,
$$V_\lambda (p)=V_\lambda (0)\mathrm{e}^{p/p_0}.$$
(18)
Here,
$$p_0=\frac{k_BT}{2L\sigma _{\text{eff}}}$$
(19)
is the characteristic decoherence constant (โdecoherence pressureโ), where $`L`$ denotes the distance between the gratings and $`\sigma _{\text{eff}}`$ corresponds to the effective cross section Hackermรผller et al. (2003b). This pressure-dependent decay of the visibility has indeed been confirmed experimentally for C<sub>70</sub> molecules Hackermรผller et al. (2003b); Hornberger et al. (2003), in excellent agreement with the theoretical predictions (Fig. 7).
Studies of collision-induced decoherence in a Talbot-Lau interferometer not only represent an outstanding method to observe the gradual disappearance of quantum-interference effects while having full control over both the source and the strength of decoherence, but also allow one to predict the environmental conditions (in this case, the maximum pressure of the surrounding gas) required to observe quantum effects for even more complex and massive objects than tested thus far. Such experiments are limited by two main factors Hornberger et al. (2003); Hackermรผller et al. (2003b). First, the velocity of the objects must be quite slow during the passage through the interferometer, in order to keep the de Broglie wavelengths long enough to allow for a sufficient degree of diffraction by practically realizable gratings. Second, the pressure $`p`$ of the residual gas in the interferometer must be low enough to maintain sufficient visibility of the interference pattern, i.e., we must have $`O(p)=p_0`$, see Eq. (19). Since both limits are purely technical and can be precisely quantified, there is no indication for any fixed quantum-classical boundary in this case other than the observational limit determined by environmental decoherence, for which rigorous theoretical estimates can be given. Decoherence allows for an exact specification of where the quantum-to-classical transition occurs and what needs to be done to move the boundary.
In fact, the envelope for the observation of the wave nature of mesoscopic molecules has recently been pushed even further in experiments demonstrating quantum interference fringes for the important biomolecule tetraphenylporphyrin C<sub>44</sub>H<sub>30</sub>N<sub>4</sub> (with mass $`m=614`$ amu and a width over 2 nm) and for the fluorinated fullerene C<sub>60</sub>F<sub>48</sub> (mass $`m=1632`$ amu, 108 atoms) Hackermรผller et al. (2003a). While tetraphenylporphyrin is the first-ever biomolecule whose wave nature has been demonstrated experimentally, fluorofullerenes are the most massive and complex molecules to exhibit quantum behavior thus far. Theoretical estimates for the maximum residual gas pressure that would still allow for the observation of interference fringes for even larger biological objects, up to the size of a rhinovirus, have been given by Hackermรผller *et al.* Arndt et al. (2002); Hackermรผller et al. (2003b) (see Fig. 9) and appear to be realizable even with the currently available technology in Talbot-Lau interferometry Hornberger et al. (2003); Hackermรผller et al. (2003b). One might extrapolate even further and speculate about the feasibility of interference experiments involving human cells, with an average weight and size on the order of $`10^{15}`$ amu and $`10^4`$ nm, respectively. While this is certainly beyond the existing technology, there is no reason to assume that such experiments should be impossible.
#### II.3.3 Implications of the C<sub>70</sub> interference experiments
The described matter-wave interferometry experiments have led to three crucial results:
1. Interference patterns are observed for particles that clearly reside in the โlump of matterโ category.
2. These patterns are due to single-particle (rather than interparticle) interference effects.
3. Any observed disappearance (or absence) of interference patterns can be well understood as resulting from decoherence and can be explicitly controlled and quantified.
Thus there is no theoretical or experimental indication for any fundamental limit on the ability of objects to exhibit quantum behavior (i.e., a wave nature) if these objects are sufficiently shielded from the decohering influence of their environment. Result (2) shows that the initial wave function describing the individual molecule evolves into a spatially extended wave function after passage through the diffraction grating, namely, into a superposition of โclassicalโ localized position states that each correspond to the molecule being in a specific region of space. The gradual disappearance of interference due to controlled interaction with the environment can be understood as entanglement between the different relative states of the environment and the individual components $`|๐ฑ_{\text{C}70}`$ in the superposition. It is important to note that all components $`|๐ฑ_{\text{C}70}`$ are still present regardless of the environmental interaction โ decoherence is in principle fully reversible, as experiments on coherent state-vector revival have shown (see, e.g., Raimond et al. (1997)).
### II.4 Bose-Einstein condensation
As a third example, we shall discuss Bose-Einstein condensation (BEC). While this effect had been predicted theoretically already in the 1920s by Einstein Einstein (1924, 1925a, 1925b) based on ideas by Bose Bose (1924), explicit experimental verification succeeded only in 1995 Bradley et al. (1995); Davis et al. (1995); Anderson et al. (1995); Bradley et al. (1997). When an atomic bosonic gas confined by a magnetic trap is cooled down to very low temperatures, the de Broglie wavelength $`\lambda _{\text{dB}}=(2\pi \mathrm{}^2/mk_BT)^{1/2}`$ associated with each atom becomes long in comparison with the interparticle separation. At a precise temperature in the $`100`$ nK range, the collection of atoms can undergo a quantum-mechanical phase transition to a condensate in which the atoms lose their individuality and all occupy the same quantum state. Then a macrosopic number of atomsโlarge condensates can contain of the order of $`10^7`$ atomsโis described by a single $`N`$-particle wave function with a phase,
$$\mathrm{\Psi }_N(๐ซ_1,๐ซ_1,\mathrm{},๐ซ_N)=\mathrm{e}^{\mathrm{i}\mathrm{\Phi }}\underset{i=1}{\overset{N}{}}|\psi (๐ซ_i)|,$$
(20)
i.e., as a product of $`N`$ identical single-particle wave functions $`\psi (๐ซ)`$. As a consequence, BECs can directly exhibit quantum behavior. For instance, two condensates released from adjacent traps can overlap and form a gas-density interference pattern due to the phase difference between the two wave functions (Fig. 10) Javanainen and Yoo (1996); Andrews et al. (1997); Rรถhrl et al. (1997); Javanainen (2005); Saba et al. (2005). Recently, Bose-Einstein โdouble-slitโ interferometers have been experimentally realized Shin et al. (2004) and theoretically analyzed Collins et al. (2005). Here, a single condensate is coherently split (corresponding to the diffraction stage in the double-slit experiment) and then allowed to recombine, which leads to the observation of interference fringes (Fig. 10).
#### II.4.1 Macrosopic number-difference superpositions using Bose-Einstein condensates
Various methods have been proposed for the creation of BEC-based Schrรถdinger cat states in form of a superposition of states with macroscopically distinguishable numbers of particles Cirac et al. (1998); Ruostekoski et al. (1998); Gordon and Savage (1999); Dunningham and Burnett (2001); Calsamiglia et al. (2001); Louis et al. (2001); Micheli et al. (2003). BECs are particularly suitable for the generation and the study of Schrรถdinger cat states, for several reasons. First, as BECs involve up to $`10^7`$ atoms, such superpositions would be the most macrosopic ones ever observed. Second, the condensate is described by a single coherent wave function that pertains to a controllable number of atoms and possesses an extremely long coherence time (up to 10โ20 s). Third, the sources of decoherence (mostly loss of particles from the condensate) are fairly well-understood and potentially sufficiently controllable through suitable environmental engineering and trap design Ruostekoski and Walls (1998); Kuang et al. (1999); Dalvit et al. (2000).
The typically suggested scheme to create superpositions of macrosopically distinguishable states using BECs involves the creation and manipulation of interacting two-species condensates, i.e., BECs in which the atoms possess two different internal states $`|A`$ and $`|B`$. Experimental realizations of two-species BECs often employ the two hyperfine sublevels $`|F,m_F=|2,1`$ and $`|1,1`$ of <sup>87</sup>Rb. The early proposal by Cirac *et al.* Cirac et al. (1998) (similiar models have been suggested, for example, in Ruostekoski et al. (1998); Gordon and Savage (1999); Dunningham and Burnett (2001); Micheli et al. (2003); Jack and Yamashita (2005)) is based on a Josephson-like coupling between the two species that leads to a number-difference superposition of the form
$$|\mathrm{\Psi }=\frac{1}{\sqrt{2}}\left(|n_A,Nn_A+\mathrm{e}^{\mathrm{i}\phi }|Nn_A,n_A\right),$$
(21)
where $`|n_A,n_B`$ is the occupation-number state representing $`n_A`$ atoms of type A and $`n_B`$ atoms of type B, and $`N=n_A+n_B`$ is the total number of atoms. This represents a superposition of two states which differ by a macrosopic number $`|N2n_A|`$ of atoms of a certain type (A or B). Then $`n_A=0`$ or $`n_A=N`$ would correspond to a maximally entangled $`N`$-particle GHZ-type state Greenberger et al. (1990) and thus the most โcat-likeโ state
$$|\mathrm{\Psi }=\frac{1}{\sqrt{2}}\left(|N,0+\mathrm{e}^{\mathrm{i}\phi }|0,N\right).$$
(22)
Another scheme for the creation of macrosopic BEC superpositions that uses a single-component BEC in a double well (with possible generalizations to $`M`$ wells) has been described in Mahmud et al. (2004, 2005) (see also Polkovnikov et al. (2002); Polkovnikov (2003)). Here, a laser-induced phase shift is imprinted on the condensate in one of the wells, followed by a change of barrier height. This is predicted to lead to a superposition of the form
$$|\mathrm{\Psi }=\frac{1}{\sqrt{2}}\left(|n_L,Nn_L+\mathrm{e}^{\mathrm{i}\phi }|Nn_L,n_L\right),$$
(23)
where $`|n_L,n_R`$ is the number state corresponding to $`n_L`$ ($`n_R`$) atoms in the left (right) well. Again, $`n_L`$ determines the degree of entanglement, with $`n_L=0`$ or $`n_L=N`$ corresponding to maximal โcatness.โ Even the possibility of creating a coherent superposition of a macroscopic number of atoms with a macroscopic number of molecules using photoassociation in BECs (i.e., the absorption of a photon by two atoms, leading to the formation of a two-atomic bound molecule) has been indicated Calsamiglia et al. (2001).
To detect a BEC cat state, one might in principle envision experiments similiar to those measuring GHZ spin states Mermin (1990); Dalvit et al. (2000); Sackett et al. (2000), although this would be very difficult to carry out in practice for the larger values of $`N`$ relevant to BEC superpositions. Instead, as pointed out in Dalvit et al. (2000), one could first confirm that measurement statistics indeed give equal likelihoods for the two cat-state terms $`|N,0`$ or $`|0,N`$. If the system can also be observed to (approximately) return to its initial state after unitary evolution over a period that is an integer multiple of the time needed for the generation of the cat state, this would provide strong indications for the presence of a cat state.
#### II.4.2 Decoherence of BEC superpositions
To date, Schrรถdinger-cat states using BECs have not been realized experimentally, although much progress has been made (see, for example, Albiez et al. (2004)). Dissipation and decoherence effects are still too strong to allow for a direct observation of superpositions and will continue to constitute the dominant limit on the size of number-difference Schrรถdinger cats. These environmental effects are mainly due to elastic and inelastic scattering between condensate and noncondensate atoms.
Elastic collisions with noncondensate atoms under conservation of the number of condensate atoms lead to phase damping and thus to the destruction of the coherent superposition. The reduced density matrix in the number basis then decoheres according to Louis et al. (2001)
$$m|\widehat{\rho }(t)|n=\mathrm{e}^{(mn)^2\kappa t}m|\widehat{\rho }(0)|n\mathrm{e}^{\mathrm{i}\omega (mn)t},$$
(24)
i.e., the off-diagonal elements $`mn`$ decay with a decoherence rate that scales with the square of the number difference, $`(mn)^2`$.
Furthermore, inelastic collisions with noncondensate atoms lead to a loss of atoms from the condensate, which diminishes coherence. Again, the larger the number difference in the superposition $`(|n,Nn+|Nn,n)/\sqrt{2}`$ is (i.e., the closer $`n`$ is to $`0`$ or $`N`$), the more sensitive the state is to atom loss (see, for example, the detailed analysis in Dunningham and Burnett (2001)). In the limit of the maximally entangled state $`\left(|N,0+|0,N\right)/\sqrt{2}`$, already the loss of a single atom of, say, type 1 completely destroys the coherent superposition, since
$$\widehat{a}_1\left(|N,0+\mathrm{e}^{\mathrm{i}\phi }|0,N\right)/\sqrt{2}=\sqrt{N/2}|N1,0,$$
(25)
where $`\widehat{a}_1`$ is the destruction operator for particles of type 1.
Thus both decoherence effects will usually limit the size $`N`$ (i.e., the number difference) of superpositions of the form $`\left(|N,0+|0,N\right)/\sqrt{2}`$. In a detailed analysis that combines the two forms of scattering processes, Dalvit *et al.* Dalvit et al. (2000) have estimated the decoherence rate $`\tau _d^1`$ for an optimal number-difference superposition $`\left(|N,0+|0,N\right)/\sqrt{2}`$ in a standard harmonic trap due to a โthermal cloudโ of $`N_{\text{nc}}`$ noncondensate atoms as
$$\tau _d^1a^2N_{\text{nc}}N^2,$$
(26)
where $`a`$ is the scattering length. This leads to very short decoherence times even for moderate environment and condensate sizes Dalvit et al. (2000); Louis et al. (2001). For example, for $`N_{\text{nc}}=10`$ and $`N=10^3`$, $`\tau _d`$ is of the order of milliseconds. For larger Schrรถdinger cats with $`N=10^7`$ and a thermal cloud containing $`N=10^4`$ noncondensed atoms, $`\tau _d10^{13}`$ s.
However, several schemes exist to significantly reduce the decoherence rate and to thus render it quite likely that BEC-based number-difference Schrรถdinger cat states could indeed be observed in future experiments; for example:
1. The construction of modified traps that allow for a faster evaporation of the thermal cloud Dalvit et al. (2000).
2. Generation of number-difference cat states via the creation of macrosopic superpositions of relative-phase states that are not only much less sensitive to atom loss, but might even require such loss Dunningham and Burnett (2001).
3. A โsymmetrizationโ of the environment to reduce the effective size of the thermal cloud Dalvit et al. (2000).
4. Sufficiently fast generation of the cat state Micheli et al. (2003).
The key lesson to be learned from the example of BEC-based Schrรถdinger-cat states is that, nonwithstanding the fact that such superpositions have not (yet) been explicitly documented in experiments, the physics of these states and the required conditions to create them is very well understood. The failure to experimentally generate these states with currently available setups is well-explained by decoherence models that provide precise numerical estimates for the type of experimental arrangements and parameter ranges that would be required to observe Schrรถdinger-cat states using BECs. Similiar to the case of studying the feasibility of matter-wave interferometry with larger molecules than those investigated thus far (see Sec. II.3), decoherence is the key tool for a precise prediction of the physical conditions required for the experimental observation of superpositions of macrosopically distinct states.
### II.5 Analysis of the degree of macroscopicity of the experimentally achieved superpositions
In the following, let us compare the degree of macrosopic distinctness of the states in the superpositions encountered in the experiments with SQUIDs, diffracted molecules, and BECs. We will use the combination of the two measures introduced in Sec. II.1, namely, the difference $`๐ฎ_{\text{ext}}`$ in a relevant extensive quantity between the states in the superposition relative to an appropriate microsopic reference value, and the degree of entanglement $`๐ฎ_{\text{ent}}`$ present in the multi-particle superposition.
For the SQUID experiments (Sec. II.2), choosing the total magnetic moment to be the relevant extensive variable, the two states $`|L`$ and $`|R`$ differ by about $`10^{10}\mu _B`$ in the experiment by Friedman *et al.* Friedman et al. (2000). Taking the Bohr magneton $`\mu _B`$ as the reference unit, the extensive difference $`๐ฎ_{\text{ext}}`$ between the two states is thus of the order of $`10^{10}`$. The degree of entanglement $`๐ฎ_{\text{ent}}`$ in the multi-Cooper-pair state can be estimated to be of the order of the number of Cooper pairs, i.e., $`10^9`$.
In the case of diffraction of C<sub>70</sub> molecules (Sec. II.3), a suitable extensive quantity would be the center-of-mass displacement between the two paths through the interferometer, which we can estimate to be on the order of 1 mm (corresponding to the lateral width of the molecular beam Hackermรผller et al. (2003b)) relative to the size of the molecule of about 1 nm, which yields a value for $`๐ฎ_{\text{ext}}`$ on the order of $`10^6`$. The degree of entanglement $`๐ฎ_{\text{ent}}`$ is essentially given by the number of microsopic constituents in the molecule, $`3\times 6\times 7010^3`$.
For BEC two-species superpositions that use the two hyperfine sublevels $`|F,m_F=|2,1`$ and $`|1,1`$ of <sup>87</sup>Rb atoms (Sec. II.4), a suitable extensive variable would be the total difference in angular momentum due to the hyperfine splitting, in units of $`\mathrm{}`$, which is on the order of the number $`N`$ of atoms in the condensate, which can be as large as $`10^7`$. Thus the maximum $`๐ฎ_{\text{ext}}`$ is on the order of $`10^7`$. The degree of entanglement $`๐ฎ_{\text{ent}}`$ is again suitably measured by the number of nucleons and electrons in the condensate, which is of the order of $`100N`$ for <sup>87</sup>Rb. Note, however, that such superpositions have not yet been experimentally achieved.
All values are summarized in Table 1. We see that the SQUID experiments allow for superpositions that are about ten orders of magnitude โmore macrosopicโ (in the sense defined above) than those achieved by molecular interferometry. On the other hand, the latter experiments lead to a direct realization of spatial superpositions, which are often considered to be more โcounterintuitiveโ than the superposition of superconducting currents, since position appears to be the dominant definite quantity in our observation of the macroworld. The ubiquitous perception of definiteness in position space has even led some to postulate a fundamentally preferred role to position. For example, Bell Bell (1982) stated that โin physics the only observations we must consider are position observations, if only the positions of instrument pointers.โ A similiar idea underlies the spatial localization mechanism in the GRW theory and is reflected in the concept of definite particle trajectories in Bohmian mechanics.
Superpositions involving two-species BECs, if experimentally realized, would come close to the degree of macroscopicity achieved in SQUIDs. This result can be understood by noting the striking analogies between the two experiments. In both cases, the multi-particle system (the superconducting material in SQUIDs, or the atomic gas in BECs) is cooled down to extremely low temperatures near absolute zero. The two macrosopically distinguishable states (currents of opposite direction in SQUIDs, or different atom species in BECs) are coupled by a classically impenetrable barrier of the Josephson-junction type. In both experiments, this essentially leads to Schrรถdinger-cat states of the form
$$|\mathrm{\Psi }=\frac{1}{\sqrt{2}}\left(|N,0+\mathrm{e}^{\mathrm{i}\phi }|0,N\right),$$
(27)
where the number state $`|N,0`$ denotes $`N`$ particles (Cooper pairs in SQUIDs, or atoms in BECs) being in the first macroscopically distinguishable state (representing a clockwise current in the SQUID, or the hyperfine sublevel $`|F,m_F=|2,1`$ in BECs), and no particles being in the second state (corresponding to a counterclockwise current in the SQUID, or the hyperfine sublevel $`|F,m_F=|1,1`$ in BECs).
## III The status of physical collapse models
All existing interpretations of quantum mechanics can be viewed as either adding formal rules<sup>1</sup><sup>1</sup>1As, for example, done in the Copenhagen interpretation (that formally postulates a collapse, but regards it merely as an โincrease of information,โ rather than as a physical process, since it interprets the wave function as representing a probability amplitude), Bohmian mechanics, modal interpretations, and consistent-histories interpretations. or physical elements (as in collapse models) to the axioms of minimal quantum theory stated in the Introduction. With respect to the โformalโ category, if the minimal theory can be shown to be sufficient to explain and predict all our observations, there is clearly no empirical reason for introducing purely formal additives. While a similiar argument can be made regarding the โphysicalโ category, collapse theories might lead to observable deviations from Schrรถdinger dynamics and could thus be experimentally tested. In both cases, of course, there may be conceptual reasons that motivate the added elements, for example a desire to resolve a felt โweirdnessโ in the existing quantum theory. While we respect this motivation, we hope to show that in fact the minimal theory is sufficient to resolve the problems without requiring any such additions.
The increasing size of physical systems for which interference effects have been observed imposes bounds on the parameters used in collapse models. However, the current experiments demonstrating mesoscopic and macrosopic interference are still quite far away from disproving the existing collapse theories. For example, even the C<sub>70</sub> diffraction experiments described in Sec. II still fall short of ruling out continuous spontaneous localization models Pearle (1989); Diรณsi (1989); Ghirardi et al. (1990) (which lead to the strongest deviations from Schrรถdinger dynamics among all physical collapse theories) by eleven orders of magnitude Adler (2004). A recently proposed mirror-superposition experiment by Marshall *et al.* Marshall et al. (2003) that might lead to a superposition involving of the order of $`10^{14}`$ atoms still fails to rule out continuous spontaneous localization models by about six orders of magnitude Bassi et al. (2005). The superpositions observed in coherent quantum tunneling in SQUIDs also appear to be compatible with dynamical reduction models, since the spatial localization mechanism would only result in a small reduction of the supercurrent below the detectable level due to a breaking-up of Cooper pairs, but not in an approximate reduction onto one of the current states Rae (1990); Buffa et al. (1995); Bassi and Ghirardi (2003). However, given the rapid development of experiments that propose to demonstrate quantum superpositions on increasingly large scales, it appears to be only a matter of time to probe the range relevant to a test of physical reduction models.
It is important to note that no deviations from linear Schrรถdinger dynamics have ever been observed that could not also be explained (at least in principle) as apparent deviations due to decoherence. In fact, it would be very difficult to distinguish collapse effects from decoherence, since the large number of atoms required for the collapse mechanism to be effective also leads to strong decoherence Tegmark (1993); Benatti et al. (1995); Bassi and Ghirardi (2003). It would therefore be necessary to isolate the system of interest extremely extremely well from its environment, such that decoherence effects can be neglected with respect to the environment-independent localization mechanism. Even in this case it might be difficult to exclude the influence of decoherence due to, for example, thermal emission of radiation, as demonstrated in the case of fullerene and C<sub>70</sub> interferometry Hackermรผller et al. (2004); Hornberger et al. (2005).
This leaves physical collapse theories, at least so far, in the speculative realm, with the added difficulty of obtaining relativistic generalizations Bassi and Ghirardi (2003). Certainly, such collapse mechanisms might be discovered in the future. However, in the absence of positive experimental evidence for such effects, and given the viable option of constructing a quantum theory consistent with all observations from the minimal formalism alone (a strategy advocated in this paper), the need for a postulated collapse effect, with free parameters tuned such as to avoid inconsistencies with the observation (or nonobservation) of superpositions, appears rather doubtful.
## IV Emergence of probabilities in a relative-state framework
The question of the origin and meaning of probabilities in a relative stateโtype interpretation that is based solely on a deterministically evolving global quantum state, and the problem of how to consistently derive Bornโs rule in such a framework, has been the subject of much discussion and criticism aimed at this type of interpretation (see, e.g., Kent (1990)). Several decoherence-unrelated proposals have been put forward in the past to elucidate the meaning of probabilities and to arrive at the Born rule in an explicit or implicit relative-state context (see, for instance, Everett (1957); Hartle (1968); DeWitt (1971); Graham (1973); Geroch (1984); Deutsch (1999)). However, it is highly controversial whether these approaches are successful and represent a noncircular derivation Stein (1984); Kent (1990); Squires (1990). A derivation that is only based on the non-probabilistic axioms of quantum mechanics and on elements of classical decision theory has been presented by Deutsch Deutsch (1999). It was critized by Barnum *et al.* Barnum et al. (2000), but was subsequently defended by other authors Gill (2003); Wallace (2003b) and embedded into an operational framework by Saunders Saunders (2002). It is fair to say that no decisive conclusion appears to have been reached as to the success of these derivations.
Initially, decoherence was thought to provide a natural account of the probability concept in a relative-state framework, by relating the diagonal elements of the decohered reduced density matrix to a collection of possible โeventsโ that can be reidentified over time, and by interpreting the corresponding coefficients as relative frequencies of branches, thus leading to an interpretation of quantum probabilities as empirical frequencies Zurek (1998); Deutsch (1999). However, as it has been pointed out before Zeh (1997); Zurek (2003a); Schlosshauer (2004), this argument cannot yield a noncircular derivation of the Born rule, since the formalism (in particular, the trace operation) and interpretation of reduced density matrices presume this rule.
The solution to the problem of understanding the meaning of probabilities and of deriving Bornโs rule in a relative-state framework must therefore be sought on a much more fundamental level of quantum mechanics. Since this framework presumes nothing besides the unitarily evolving state vector itself, the solution should preferably be derived solely from properties of this quantum state. However, while we would like to assign probabilities to โoutcomes of measurementsโ on a local system (i.e., probabilities for the system to be found in a certain state), the global quantum state usually contains a high degree of environmental entanglement, i.e., there exists no state vector that could be assigned to the local system alone. Still, we obviously talk regularly of the โstate of the system,โ and we must therefore distinguish this notion of state from the quantum state vector itself. Following the relative-state viewpoint, the local โeventsโ of the system (or its possible โstates of the systemโ) are then typically identified with the relative-state components of the global state vector in the Hilbert subspace corresponding to the system.
The recent enormous advances in the field of quantum information theory, especially in the understanding of the properties and implications of quantum entanglement, have shed some light on how one might proceed. Quantum information theory has established the notion that quantum theory can be viewed as a description of what, and how much, โinformationโ Nature is willing to proliferate. For example, a peculiar feature of quantum mechanics is that complete knowledge of the global pure bipartite quantum state $`|\mathrm{\Psi }=\left(|\alpha _1|\beta _1+|\alpha _2|\beta _2\right)/\sqrt{2}`$ itself does not appear to contain information about the โabsoluteโ state of one of the subsystems. This hints at ways how a concept of โignorance,โ and therefore of probability, may emerge directly from the quantum feature of entanglement without any classical counterpart.
This idea has recently been developed in a series of papers by Zurek Zurek (2003a, b, 2004b, 2004a), leading to a proposal for a derivation of Bornโs rule. As pointed out by the present author Schlosshauer and Fine (2005); Schlosshauer (2004) and made more explicit in the most recent of Zurekโs articles on this topic Zurek (2004a), the derivation is still based on certain assumptions that are not contained in the basic measurement-free relative-state framework of quantum mechanics. One might argue how strong these assumptions are. Zurek himself, for example, considers some of them to be โfactsโ and regards others as โnaturalโ and โmodestโ Zurek (2004a); a somewhat more critical position with respect to some of the assumptions has been assumed by the present Schlosshauer and Fine (2005) and other authors Barnum (2003); Mohrhoff (2004). Granted these assumptions, however, we consider Zurekโs proposal a very promising approach towards a deeper understanding of the origin of quantum probabilities, and we shall therefore outline the basic ideas and assumptions in the following (a more detailed description and discussion of the approach can be found in Zurek (2003a); Schlosshauer and Fine (2005); Schlosshauer (2004); Zurek (2004a)).
Zurekโs derivation is based on a particular symmetry property (referred to as โenvironment-assisted envariance,โ or โenvarianceโ for short) of composite quantum states, which is used to infer complete ignorance about the state of the subsystem. The derivation relies on a study of the properties of a composite entangled state and therefore intrinsically requires the decomposition of the Hilbert space into subsystems and the usual tensor-product structure. The core result to be established is the following. Given a bipartite product Hilbert space $`H_{๐ฎ_1}H_{๐ฎ_2}`$ and a completely known composite pure state in the diagonal Schmidt decomposition
$$|\mathrm{\Psi }=\left(e^{i\phi _1}|\alpha _1_1|\beta _1_2+e^{i\phi _2}|\alpha _2_1|\beta _2_2\right)/\sqrt{2},$$
(28)
where the $`|\alpha _i_1`$ and $`|\beta _i_2`$ are orthonormal basis vectors that span the Hilbert spaces $`H_{๐ฎ_1}`$ and $`H_{๐ฎ_2}`$, the probabilities of obtaining either one of the relative states $`|\alpha _1_1`$ and $`|\alpha _2_1`$ (identified with the โeventsโ of interest to which probabilities are to be assigned (Zurek, 2004b, p. 12); see also Schlosshauer and Fine (2005)) are equal. Given this result, generalizations to higher-dimensional Hilbert spaces and to the case of unequal absolute values of the Schmidt coefficients in Eq. (28) can be achieved in a rather straightforward way Zurek (2004a).
This result is established in two key steps. First, a few simple assumptions (Zurekโs โfactsโ Zurek (2004a)) are introduced that connect the global quantum state $`|\mathrm{\Psi }`$, Eq. (28), to the โstate of the systemโ $`๐ฎ_1`$. This is necessary because, as mentioned above, the global quantum state is all that the pure state-vector formalism of quantum mechanics provides for the description of a bipartite system containing entanglement. The following assumptions are made about the โstate of the systemโ $`๐ฎ_1`$. First, this state is completely determined by the global quantum state, Eq. (28); second, it specifies all measurable properties of $`๐ฎ_1`$, including probabilities of outcomes of measurements on $`๐ฎ_1`$; and third, unitary transformations can change it only if they act on $`๐ฎ_1`$ (see Schlosshauer and Fine (2005) for a discussion of this last assumption).
Granted these three assumptions, one can show that measurable properties of $`๐ฎ_1`$ can depend neither
1. on the phases $`\phi _i`$ in Eq. (28), such that we can assume the simplified form
$$|\mathrm{\Psi }=\left(|\alpha _1_1|\beta _1_2+|\alpha _2_1|\beta _2_2\right)/\sqrt{2}$$
(29)
for our purpose of discussing probabilities associated with $`๐ฎ_1`$;
2. nor on whether $`|\alpha _1_1`$ is paired with $`|\beta _1_2`$ or $`|\beta _2_2`$, i.e., the unitary transformation acting on $`๐ฎ_1`$ that changes the quantum state vector
$$|\mathrm{\Psi }=\left(|\alpha _1_1|\beta _1_2+|\alpha _2_1|\beta _2_2\right)/\sqrt{2}$$
(30)
into
$$|\mathrm{\Psi }^{}=\left(|\alpha _2_1|\beta _1_2+|\alpha _1_1|\beta _2_2\right)/\sqrt{2}$$
(31)
cannot have altered the state of $`๐ฎ_1`$.
In a way, result (2) already indicates a feature of ignorance about the state of $`๐ฎ_1`$, since interchanging the potential โoutcomesโ $`|\alpha _i_1`$ through local operations performed on $`๐ฎ_1`$ does not change any measurable properties of $`๐ฎ_1`$ and can therefore be viewed as leading to a form of โobjective indifferenceโ among these outcomes. It is important to note that this effect is crucially dependent on the feature of entanglement. In a nonentangled pure state of the form $`|\mathrm{\Phi }=\left(|\varphi _1+e^{i\phi }|\varphi _2\right)\sqrt{2}`$, the phase $`\phi `$ must of course not be ignored (and would be measurable in a suitable interference experiment), and therefore the system described by the โswappedโ state $`|\mathrm{\Phi }^{}=\left(|\varphi _2+e^{i\phi }|\varphi _1\right)\sqrt{2}`$ is clearly physically different from that represented by the original state $`|\mathrm{\Phi }`$.
To make the above argument more precise, in the second key step of the derivation, the notion of probabilities of the outcomes $`|\alpha _i_1`$ in a measurement performed on $`๐ฎ_1`$ (previously only subsummed under the general heading โmeasurable properties of $`๐ฎ_1`$โ) is now explicitly connected to the global state vector via an additional assumption. In Zurek (2004a), Zurek offers three possible choices for this assumption, of which we should quote one (see also Barnum (2003)). Namely, it is assumed that the form of the Schmidt product states $`|\alpha _i_1|\beta _i_2`$ appearing in Eq. (28) implies that the probabilities for $`|\alpha _i_1`$ and $`|\beta _i_2`$ must be equal. Given this assumption and using result (2) above, it can be readily established Schlosshauer and Fine (2005); Schlosshauer (2004); Zurek (2004a) that the probabilities for $`|\alpha _1_1`$ and $`|\alpha _2_1`$ must be equal, thus completing the derivation.
As we have pointed out elsewhere Schlosshauer and Fine (2005), the need for the final assumption may be considered a reflection of the well-worn phrase that a transition from a nonprobabilistic theory (such as quantum mechanics solely based on deterministically evolving state vectors) to a probabilistic theory (that refers to โprobabilities of outcomes of local measurementsโ) requires, at some stage, to โput probabilities in to get probabilities out.โ However, in the quantum setting, this introduction of a probability concept has a far more objective character than in the classical case. While in the latter setting probabilities refer to subjective ignorance in spite of the existence of a well-defined state (see also Sec. V), in the quantum case all that is available, namely, the global entangled quantum state, is perfectly known. The objectivity of ignorance in quantum mechanics can thus be viewed as a consequence of a form of โcomplementarityโ between local and global observables Zurek (2004a) and could help explain the fundamental need for a probabilistic description in the quantum setting despite the deterministic evolution of the global state vector.
It is the great merit of Zurekโs proposal to have emphasized this objective character of quantum probabilities arising from the feature of quantum entanglement. While the precise role and importance of the assumptions entering the derivation as well as the generality of the approach (given, e.g., the focus on Schmidt decompositions) would benefit from further discussion and analysis, the approach definitely sheds an interesting and new light on the nature of quantum probabilities.
## V Objectification of observables in a relative-state framework
A characteristic feature of classical physics is the fact that the state of a system can be found out and agreed upon by many independent observers (with all of them initially completely ignorant about the state) without disturbing this state. In this sense, classical states preexist objectively, resulting in our notion of โclassical reality.โ In contrast, as is well known, measurements on a closed quantum system will in general alter its stateโunless, of course, the observer chooses to measure, by pure luck or prior knowledge, an observable with an eigenstate that coincides with the state of the system. It is therefore impossible to regard quantum states of a closed system as existing in the way that classical states do. This raises the question of how classical reality emerges from within the quantum substrate, i.e., how observables are โobjectifiedโ in the above sense.
In a first step, the decoherence program, in particular the stability criterion and the more general formalism of the โpredictability sieveโ Zurek (1981, 1982, 1993, 1998, 2003a); Schlosshauer (2004) (see also Sec. II.2.4), has provided an answer to the question of why only a certain subset of the possible states in the Hilbert space of the system are actually observed. Taking into account the openness of the system and the form of the system-environment interaction is crucial in determining a set of preferred stable states of the system. This supplies an elegant and physically motivated solution to the problem of the preferred basis, an issue that has often been used to challenge the feasibility of relative-state interpretations Kent (1990); Stapp (2002). Nonetheless, the problem sketched in the previous paragraph remains, as any direct measurement performed on the system would, in general, still alter the state of the system.
The important next step is therefore to realize that in most (if not all) cases observers gather information about the state of a system through indirect observations, namely, by intercepting fragments of environmental degrees of freedom that have interacted with the system in the past and thus carry information about the state of the system Zurek (1993, 1998, 2000, 2003a). Probably the most common example for such an indirect acquisition of information is the visual registration of photons that have scattered off from the object of interest (see also Sec. VI.3). Similiar to the case of decoherence, the recognition of the openness of quantum systems is therefore crucial. However, the role of the environment is now broadened, namely, from the selection of preferred states for the system of interest and the dislocalization of local phase coherence, to the transmission of information about the state of the system. The idea is then to show how, and which, information is both redundantly and robustly stored in a large number of distinct fragments of the environment in such a way that multiple observers can retrieve this information without disturbing the state of the system, thereby achieving effective classicality of the state.
This approach has recently been developed under the labels of โenvironment as a witnessโ (i.e., the recognition of the role of the environment as a communication channel) and โquantum Darwinismโ (namely, the study of what information about the system can be stably stored and proliferated by the enviroment) Zurek (1993, 1998, 2003a, 2004b); Ollivier et al. (2004a, b); Blume-Kohout and Zurek ; Blume-Kohout and Zurek (2005). To explicitly quantify the degree of completeness and redundancy of information imprinted on the enviroment, the measure of (classical Ollivier et al. (2004a, b) or quantum Zurek (2003a); Blume-Kohout and Zurek ; Blume-Kohout and Zurek (2005)) mutual information has usually been used. Roughly speaking, this quantity represents the amount of information (expressed in terms of Shannon or von Neumann entropies) about the system $`๐ฎ`$ that can be acquired by measuring (a fragment of) the environment $``$. Note that the amount of information contained in each fragment is always somewhat less Blume-Kohout and Zurek (2005) than the maximum information provided by the system itself (as given by the von Neumann entropy of the system).
The measure of classical mutual information is based on the choice of particular observables of $`๐ฎ`$ and $``$ and quantifies how well one can predict the outcome of a measurement of a given observable of $`๐ฎ`$ by measuring some observable on a fraction of $``$ Ollivier et al. (2004a, b). The quantum mutual information $`_{๐ฎ:}`$, used in more recent studies Zurek (2003a); Blume-Kohout and Zurek ; Blume-Kohout and Zurek (2005), can be viewed as a generalization of classical mutual information and is defined as $`_{๐ฎ:}=H(๐ฎ)+H()H(๐ฎ)`$, where $`H(\rho )=\text{Tr}(\rho \mathrm{log}\rho )`$ is the von Neumann entropy. Thus $`_{๐ฎ:}`$ measures the amount of entropy produced by destroying all correlations between $`๐ฎ`$ and $``$, i.e., it quantifies the degree of correlations between $`๐ฎ`$ and $``$. Results derived from these measures have thus far been found to be sufficiently robust with respect to the particular choice of measure Ollivier et al. (2004a, b); Zurek (2003a); Blume-Kohout and Zurek ; Blume-Kohout and Zurek (2005), although a more detailed analysis of this issue is underway Blume-Kohout and Zurek (2005).
It has been found that the observable of the system that can be imprinted most completely and redundantly in many distinct subsets of the environment coincides with the โpointerโ observable selected by the system-environment interaction (i.e., by the stability criterion of decoherence) Ollivier et al. (2004a, b); Blume-Kohout and Zurek ; Blume-Kohout and Zurek (2005). Conversely, most other states do not seem to be redundantly storable. Thus it suffices to probe a comparably very small fraction of the environment to infer a large amount of the maximum information about the pointer state of the system. On the other hand, if the observer tried to measure other observables on the same fragment, he would learn virtually nothing, as information about the corresponding observables of the system is not redundantly stored. Thus the โpointerโ states of the system play a twofold role: They are the states least perturbed by the interaction with the environment, and they are the states that can be most easily found out, without disturbing the system, by probing environmental degrees of freedom. Since the same information about the pointer observable is stored independently in many fragments of the environment, multiple observers can measure this observable on different fragments and will automatically agree on the findings. In this sense, one can ascribe (effective) objective existence to the pointer states.
The research into the objectification of observables along the lines outlined in this section is only in its beginnings. Important aspects, such as the explicit dynamical evolution of the objectification process Ollivier et al. (2004b) and the role of the assumptions and definitions in the current treatments of the โobjectification through redundancyโ idea, are currently still under investigation, as are studies involving more detailed and realistic system-environment models. However, it should have become clear that the approach of departing from the closed-system view and of describing observations as the interception of information that is redundantly and robustly stored in the environment, represents a very promising candidate for a purely quantum-mechanical account of the emergence of classical reality from the quantum domain.
## VI Decoherence in the perceptive and cognitive apparatus
If, motivated by the results of the experiments described in Sec. II, we assume the universal validity of the Schrรถdinger equation, we immediately face two related consequences:
1. We ought to reconcile this assumption with our perception of definite states in the macroworld, since now there is no underlying stochastic mechanism (of whatever nature) that would select, in an objective manner, a particular โoutcomeโ among the terms in a superposition of, say, spatially localized wave packets. There exists not only a multitude, but also interference effects between them.
2. If Schrรถdinger dynamics are universal, it is reasonable (at least from a scientifically reasonable functionalistโs standpoint) to also describe observers with their perceptive and cognitive apparatusesโincluding even what could be grouped together under the rather vague term of โconsciousnessโ von Neumann (1932); Wigner (1962); Stapp (1993); Zeh (2000)โby unitarily evolving wave functions.
Both consequences follow quite naturally from the assumption of universally exact \[consequence (1)\] and universally applicable \[consequence (2)\] Schrรถdinger dynamics. Quite generally, the preferred strategy would be to treat them jointly: Solving the โmeasurement problem,โ that is, consequence (1), posed by the assumption of a purely unitary quantum theory, by applying this very theory to the observer, i.e., consequence (2). If successful, this would lead to a โsubjectiveโ resolution of the measurement problem, i.e., to a quantum-mechanical account of why we, as observers, perceive definite states in specific bases, rather than superpositions of these states. In the opinion of this author Schlosshauer (2004) and of others (see, e.g., Zeh (1973); Zurek (1998); dโEspagnat (2000); Zeh (2000); Zurek (2003a)), this would also represent a sufficient solution to the problem.
### VI.1 General remarks
First of all, on a rather philosophical sidenote, it is clear that the familiar concepts of the world of our experience are expressed in terms of the observed specific definite states. We do not even have a concept available for what a state describing a superposition of an alive and dead cat would represent, because we have never observed such a state. While such a Schrรถdinger cat might seem exotic, we have seen that quite analogous states are realizable in the laboratory โ for example, in terms of superpositions of currents running in opposite directions in SQUIDs. As we have argued in Sec. II.2.3, the only way we can access such superpositions in terms of our concepts (and not just in mathematical terms) is through the definite current states $`|L`$ and $`|R`$ that are observable as individual preferred states of the system upon measurement.
Furthermore, it is virtually indisputable that we must describe all observations in terms of physical interactions between the observed system and the observer, i.e., by means of an appropriate interaction Hamiltonian $`\widehat{H}_{\text{int}}`$. Such interactions do not have to be, and usually are not, direct. For example, the probably most common type of observation involves the interception of a number of photons that have interacted with the object of interest in the past and whose state is thus entangled with the state of the object. These photons then contain indirect and redundantly coded information about the object that can be revealed without significantly disturbing the state of the object (see Sec. V).
If the perception of definiteness is not introduced as an extraneous postulate, but is rather understood as emerging from the unitary quantum formalism itself when observations and observers are described in physical terms, it is inevitable that attempts have to be made to analyze the cognitive apparatus in quantum-mechanical terms. It is clear that giving such an account of subjective definiteness by referring to the physical structure of observers cannot share the mathematical compactness and exactness of axiomatically introduced rules that enforce definiteness on a fundamental level of the theory. However, it is important to note that, given the paramount role of observations in quantum mechanics (mostly owing to the fact that, in general, states do not pre-exist in a classical sense), postulating such โexactโ rules is tantamount to simply avoiding a physical analysis of crucial and objective (that is, interpretation-neutral) physical processes (cf. Kentโs objections to โmany-worldsโ interpretations Kent (1990) and Wallaceโs defense Wallace (2003a)).
If a purely unitary time evolution is assumed and observations are modeled as physical interactions, the conclusion of the existence of quantum-mechanical superpositions of brain states corresponding to the different โoutcomesโ of observations is inescapable. Individual perceptions are represented by certain neuronal resting/firing patterns in the brain (see Donald (1995, 2002) for more precise definitions of this relationship). As we shall discuss in the next section, superpositions of resting and firing states of a neuron are extremely sensitive to environmental decoherence, with the resting and firing states forming the robust neuronal states. These states can thus be identified with โrecord statesโ that are capable of robustly encoding information in spite of environmental interactions Zurek (1998, 2004a). As a consequence of the practically irreversible dislocalization of phase relations between these record states through entanglement with the environment, a dynamical decoupling of these states results. This process represents an objective branching process due to physical interactions between subsystems and with the environment.
The remaining question is then how to relate this objective branching to the perceived subjective โbranches of consciousness,โ i.e., collective memory states, or โmindsโ (von Neumannโs principle of the โpsycho-physical parallelismโ von Neumann (1932)). Of course, the existence (and therefore the locality) of consciousness cannot actually be derived from the quantum-mechanical formalism. This has led some authors to conclude that the question of the relationship between subjective experience and its physical correlates can only be fully answered through the introduction of new physical laws Donald (2002). However, in the opinion of this and other authors (see, for example, Zeh (2000, 2004)), it is an entirely viable (if not compelling) strategy to postulate, within the formalism, the existence of consciousness based on the empirical fact of decohering wavefunction components in neuronal processes, by associating the robust components of the global wave function labelled by the decohered neuronal states with dynamically autonomous observers Zeh (1970, 1973); Lockwood (1996); Zurek (1998); Zeh (2000); Zurek (2003a); Zeh (2004); Zurek (2004a).
Due to the absence of more concrete theoretical and experimental insight into the physical underpinnings of the cognitive apparatus with its associated complex entities such as the โmind,โ โconsciousness,โ and even the comparably basic โrecord states,โ the above brief account of how subjective definiteness may emerge from purely unitary quantum mechanics must (at least for now) remain inherently somewhat vague and nontechnical. Fortunately, however, the main points of the argument are quite independent of, say, the precise details of the structures and dynamics of the information-processing cognitive entities, since the ubiquity and effectiveness of decoherence is likely to lead to very robust results. We shall therefore turn, in the next section, to concrete estimates for decoherence rates in neurons.
### VI.2 Decoherence of neuronal superpositions
The extremely complex network of about $`10^{11}`$ interacting neurons in the brain undoubtly comprises a major part of the cognitive machinery used for processing and storing of information obtained from sensory input. Computer models of such neuronal networks (employing a massively parallel interconnected web of โswitchesโ that are turned on and off depending on some, typically nonlinear, activation function) can exhibit rich and complex behavior similiar to that encountered in cognitive processes.<sup>2</sup><sup>2</sup>2However, as Donald Donald (2002) has pointed out, the brain should not be thought of as a deterministic classical computer with a predictable input/output pattern, since synaptic transmissions have a fairly high failure rate due to the complexity of the underlying biological processes. The large number of about $`10^{14}`$ synapses in the human brain, with each neuron firing in average several times per second, inevitably leads to a high degree of unpredictability on the โeveryday levelโ that is much more significant than effects due to pure quantum uncertainties. In particular, it is reasonable to identify the โrecord statesโ mentioned above with individual neurons or neuronal clusters. One might conjecture that ultimately all cognitive processes (and thus presumably also our perception of consciousness) are due to neuronal activity.
Thus the importance of a quantitative investigation of decoherence in neuronal states should be clear. Tegmark Tegmark (2000) has estimated decoherence rates for a superposition of a firing and non-firing neuron in the brain. The firing is represented by a large number $`N10^6`$ Tegmark (2000) of Na<sup>+</sup> ions moving across the membrane into the inside of the axon (Fig. 11). Thus, a superposition of a firing and nonfiring neuron corresponds to a spatial superposition involving $`O(N)`$ Na<sup>+</sup> ions.
The extensive difference $`๐ฎ_{\text{ext}}`$ (see Sec. II.1) can then be estimated to be on the order of $`10^2`$$`10^3`$, given by a small multiple of the thickness $`h10`$ nm of the axon membrane separating the inside and outside regions, relative to the size of a Na<sup>+</sup> ion, which is on the order of 0.1 nm. While this value for $`๐ฎ_{\text{ext}}`$ is comparably small, the degree of entanglement $`๐ฎ_{\text{ent}}`$ is somewhat closer to the values listed in Table 1. Taking it to be equal to the number of microsopic constituents, we obtain $`๐ฎ_{\text{ent}}3\times 10^7`$. Thus a neuron being in a superposition of firing and resting quite clearly falls into the macrosopic category.
The decoherence rates for this superposition as estimated by Tegmark are, as expected, extremely fast. The three main sources of decoherence in this case, namely, ion-ion scattering, ion-water collisions, and long-range Coulomb interactions due to nearby ions, all result in decoherence times on the order of $`10^{20}`$ s.
One obvious implication of fast neuronal decoherence is that coherent superpositions in neurons could never be sustained long enough to allow for some form of quantum computation. This result appears to be much more clearly established than an answer to the question of whether the relevant decoherence times are long enough to allow for quantum computation in microtubules (dynamically active structures that are a dominant part of the cytoskeleton, i.e., the internal scaffolding of cells). Suggestions in the positive, including the association of such quantum computations with the emergence of consciousness, have been put forward in Penrose (1994); Hameroff and Penrose (1996a, b), criticized in Tegmark (2000), subsequently defended in Hagan et al. (2002), and further evaluated in Rosa and Faber (2004) (see also Stapp (2000)).
However, the question much more relevant to the theme of this paper concerns the implications of neuronal decoherence for a decoherence-based account of subjective definiteness in unitary quantum mechanics โ i.e., for a subjective resolution of the โmeasurement problem.โ To this extent, let us in the following discuss a simple step-by-step quantum-mechanical account of the chain of interactions leading to the recording of a visual event in the brain.
### VI.3 Schematic sketch of the chain of interactions in visual perception and cognition
Suppose that a small number of photons interact with an object $`๐ช`$ described by a pure-state superposition of two macrosopically distinct positions. This step already can be viewed as an environmental decoherence process, where now, however, the environment assumes a crucial role as a carrier of information (see Sec. V). Due to entanglement, the combined object-photon system will be described by a superposition of the form
$$|\mathrm{\Psi }_{๐ช๐ซ}=\frac{1}{\sqrt{2}}\left(|\omega _1_๐ช|\varphi _1_๐ซ+|\omega _2_๐ช|\varphi _2_๐ซ\right),$$
(32)
where $`\omega _i`$, $`i=1,2`$, are the two distinct (small) spatial regions associated with the object, and $`|\varphi _i_๐ซ`$ denote the corresponding classically distinct collective photon states. A conceptually similiar arrangement on the mesosopic scale has explicitly been studied in experiments involving a single rubidium atom (representing the object) in a superposition of two internal levels and entangled with a cavity radiation mode (corresponding to the collection of photons) Brune et al. (1996); Raimond et al. (1997).
Initial detection of such a collection of photons in the (human) eye is associated with rhodopsin molecules in the retina. Due to their mesoscopic properties, rhodopsin molecules are subject to strong decoherence, such that already at this stage the influence of the environment will have preselected the robust states $`|\rho _i_{}`$ of the rhodopsin molecule, which correspond to certain photon detection events $`|\varphi _i_๐ซ`$. The total state $`|\mathrm{\Psi }_{๐ช๐ซ}`$ will then be given by
$$|\mathrm{\Psi }_{๐ช๐ซ}=\frac{1}{\sqrt{2}}\left(|\omega _1_๐ช|\varphi _1_๐ซ|\rho _1_{}+|\omega _2_๐ช|\varphi _2_๐ซ|\rho _2_{}\right),$$
(33)
i.e., the photon-rhodopsin interaction should lead to an (albeit, due to the influence of decoherence, very fragile) superposition of the different biochemically distinct states $`|\rho _i_{}`$ of the rhodopsin molecule.<sup>3</sup><sup>3</sup>3A search for experimental evidence for such superpositions has been suggested in Shimony (1998); for an experimental proposal, see Hilaire et al. (2002). Cf. also Thaheld (2005) for an (unconvincing) suggestion that the visual apparatus itself might trigger a physical collapse. These relative states can then be expected to be further correlated with the appropriate states $`|\nu _i_๐ฉ`$ of neuronal arrays that are mainly located in the primary visual area in the occipital lobe of the brain. Suppose that the two โeventsโ represented by the two distinct states $`|\rho _i_{}`$ of the rhodopsin molecule (corresponding to the different photon states $`|\varphi _i_๐ซ`$ that in turn carry information about the two distinct spatial regions $`\omega _i`$ of the object) are encoded by the states $`|\nu _i_๐ฉ`$, $`i=1,2`$, describing the same collection of $`N`$ neurons in two different firing/resting patterns.
As a simple example, let us take $`N=3`$, and $`|\nu _1_๐ฉ=|1_{๐ฉ_1}|0_{๐ฉ_2}|1_{๐ฉ_3}`$ and $`|\nu _2_๐ฉ=|0_{๐ฉ_1}|1_{๐ฉ_2}|1_{๐ฉ_3}`$, where $`|0_{๐ฉ_i}`$ and $`|0_{๐ฉ_i}`$ denote, respectively, the resting and firing state of the $`i`$th neuron. Then the combined state $`|\mathrm{\Psi }_{๐ช๐ซ๐ฉ}`$ will be given by
$$\begin{array}{c}|\mathrm{\Psi }_{๐ช๐ซ๐ฉ}=\frac{1}{\sqrt{2}}(|\omega _1_๐ช|\varphi _1_๐ซ|\rho _1_{}|1_{๐ฉ_1}|0_{๐ฉ_2}|1_{๐ฉ_3}\hfill \\ \hfill +|\omega _2_๐ช|\varphi _2_๐ซ|\rho _2_{}|0_{๐ฉ_1}|1_{๐ฉ_2}|1_{๐ฉ_3}).\end{array}$$
(34)
The extreme fast decoherence rate for the neurons 1 and 2 being in a superposition of firing and resting will lead to a practically irreversible dynamical decoupling of the two branches that now describe two distinct โoutcomesโ encoded by $`|\nu _i_๐ฉ`$. We may then identify these states with the basic memory states, although, strictly speaking, the physical process of actual information storage in the brain (i.e., learning) occurs only in two subsequent stages Kandel et al. (2001). First, in form of short-term memory, believed to be due to certain biochemical and electrical interactions between neurons. Second, as long-term memory that is based on actual structural changes in the brain (โneuroplasticityโ), most notably, due to the formation of new connections (synapses) between neurons and due to internal changes in the synaptic regions in individual neurons.
However, since all these processes will again be subject to strong decoherence, the essence of our argument is not altered: The states in a superposition of neuronal firing patterns will rapidly entangle with approximately orthogonal (i.e., macrosopically distinguishable) states of the environment and thus lead to the formation of locally noninterfering (that is, dynamically autonomous) branches labelled by these โoutcomeโ states. Regardless of the precise physical, chemical, biological, psychological, etc., details of perceptive and cognitive activity, it is quite clear that decoherence effects are likely to be sufficient to explain the emergence of a subjective perception of single outcomes, represented by stable, โclassical,โ record states, from a (by all accounts macrosopic) global superposition.
## VII Discussion and outlook
We have analyzed three important experimental domainsโnamely, SQUIDs, molecular diffraction, and Bose-Einstein condensationโthat have demonstrated (or at least have come very close to demonstrating) the existence of superpositions of states that can be considered macrosopically distinct in comparison with the microsopic states โtypicallyโ treated in quantum mechanics. These experiments have provided powerful examples for the validity of unitary Schrรถdinger dynamics and the superposition principle on increasingly large length scales. They have also shown how the fragility of macrosopic superpositions can be precisely understood and controlled in terms of environmental interactions and the resulting decoherence effects.
Of course, these experiments do not falsify the possibility that the Schrรถdinger equation might not be exact under all circumstances. In fact, no finite number of experiments that show the validity of unitary dynamics could ever do. To do so, a โpositive-testโ experiment would be needed that could explicitly demonstrate nonlinear deviations from the Schrรถdinger equation. Leggett Leggett (1980, 2002) has presented a Bell-type inequality that would be obeyed by what he calls the class of โmacrorealistic theories,โ while it would be violated by the predictions of purely unitary quantum mechanics. The โmacrorealisticโ class is defined to represent all theories in which macrosopic systems are always in a single definite state among a collection of possible macrosopically distinct states, and in which this definite state can be found out without perturbing the state and dynamics of the system. So one might, at least in principle, through suitable experiments be able to exclude either any such macrorealistic theory or the universal validity of the Schrรถdinger equation. Such a strategy would be similiar in spirit to the tests of Bellโs inequalities, which rule out a large class of, if not all, local realistic theories. (See Sec. 6 of Leggett (2002) and references therein for some first ideas in this direction.)
At the current stage, however, it is the opinion of the present author that, in absence of any positive evidence for deviations from unitary dynamics, combined with the continued experimental verification of increasingly large โSchrรถdinger catsโ (whose time evolution, including decoherence effects, is in perfect agreement with unitary dynamics), it appears to be not only reasonable, but moreover compelling, to entertain the possibility of a universally exact Schrรถdinger equation seriously and to fully explore the consequences of this assumption.
The experiments described in this paper have demonstrated rapid progress in achieving, controlling, and observing superpositions of increasingly distinct states. Experiments involving superpositions of classically distinguishable states of a few photons Brune et al. (1996); Raimond et al. (1997) have been followed by collective superpositions of $`10^9`$ electron pairs in SQUIDs and double slitโtype experiments using massive molecules with a large number of degrees of freedom. It is only a matter of time until number-difference superpositions involving on the order of $`10^7`$ rubidium atoms will be experimentally realized in BECs. It is rather unlikely that this progress towards experimental evidence for increasingly large superpositions will encounter any fundamental boundaries in the near future. As we have seen, the main limit seems to be given by the ability to shield the system sufficiently from the decohering influence of the environment. This limit is open to precise quantitative analysis.
In view of this situation, we may now legitimately ask what the next steps in solidifying the empirical support for a purely unitary quantum theory and its consequences might ideally look like. To this extent, we remark that superpositions of macrosopically distinct states that refer to biological (and, even more so, animate) objects seem to have been considered as particular โparadoxicalโ โ after all, Schrรถdinger chose a cat to illustrate his famous Gedanken experiment. This attitude may be traced back to several reasons. For example:
1. The โdistinctnessโ between the states referring to biological objects is usually extremely complex. Not only is the number of physical, chemical, biological, etc., differences between a dead and alive cat overwhelmingly vast, even two functionally different states of a simple biological molecule will be distinct in a large number of features. By contrast, in the examples involving inanimate objects, such as BECs and SQUIDs, the states in the superposition usually differ only in a single physical quantity, such as total angular momentum or magnetic moment.
2. While we might be willing to accept the existence of an โexoticโ superposition under extreme physical conditions (such as superconducting currents in a bulk of matter cooled down to temperatures close to absolute zero), biological objects reside in the parameter regime characteristic of the world of our everyday experience.
3. If the superposition principle is applied to human observers (specifically, superpositions of โstates of consciousness,โ etc.), we feel that our most basic intuition about possessing a unique identity has been infringed upon.
Especially in light of the first two arguments, the molecular diffraction experiments appear to be the most โnaturalโ realization of superpositions of macrosopically distinct states. In fact, as pointed out in Sec. II.3.2, interference effects have already be experimentally demonstrated for a biological molecule Hackermรผller et al. (2003a), and larger biological structures are likely to follow Arndt et al. (2002); Hackermรผller et al. (2003b) (see also Fig. 9).
However, another interesting direction could also be taken from here. As suggested for example in Shimony (1998); Leggett (2002), one might try to look for interference effects between (and thus superpositions of) biologically distinct states of the same biomolecule, rather than for the spatial superpositions demonstrated in the current molecular-diffraction setups. While such experiments would be considerably more difficult to carry out due to the required nearโin vivo environmental conditions (room temperatures, presence of a surrounding medium such as an aqueous solution, etc.), which would lead to very strong decoherence effects, there does not seem to exist a fundamental obstacle that would prevent one in principle from the realization of such superpositions in a cleverly designed setup.
Experiments that would find some basic biological structure in a superpositions of distinct states corresponding to different biological โinputsโ might in turn indicate the presence of a superposition of input signals originating from the inanimate outside world (e.g., a superposition of photon states entangled with spatially distinct states of a single object โ see Sec. VI.3). They could also provide direct empirical evidence for consequences of purely unitary quantum mechanics in the regime of more complex structures that are part of conscious (human) observers, and might therefore also ease the discomfort spelled out in item (3) above.
Given that experiments Hackermรผller et al. (2003a) have demonstrated a splitting of the localized state of a biomolecule into โbranchesโ corresponding to distinct paths, it would also be worth discussing, as Zeh Zeh (2000) puts it,
> the consequences of similar *Gedanken* experiments with objects carrying some primitive form of โcore consciousnessโ โ including an elementary awareness of their path through the slits.
In such a situation, after passage through the slits, the state of the object would be described by a superposition of spatially distinct trajectories. However, due to its awareness of the path, it would thus also be in a superposition of multiple (local) โstates of consciousness.โ Environmental scattering would then lead to entanglement with path-encoding variables (decoherence), which hence would also destroy interference effects between the โbranches of consciousness,โ and thus the different paths would be โexperiencedโ separately. In the absence of decoherence, it would be possible to coherently recombine the branches into a single localized wavepacket identical to the state before the passage through the slits. It then follows from the standard quantum-mechanical formalism that the associated object then cannot have retained any โmemoryโ of the path taken before the recombination. For related ideas using the example of neutron interferometry, see Vaidman (1998).
As it is well known, Bohr has repeatedly insisted on the fundamental role of classical concepts (see, for example, Bohr (1923, 1948)). The experimental evidence for superpositions of macrosopically distinct states on increasingly large length scales counters such a dictum. Superpositions appear to be novel and individually existing states, often without any classical counterparts. Only the physical interactions between systems then determine a particular decomposition into classical states from the view of each particular system. Thus classical concepts are to be understood as locally emergent in a relative-state sense and should no longer claim a fundamental role in the physical theory.
We have already widely acknowledged, based on experimental evidence, the fundamental nonlocality of the quantum world, in spite of the utterly nonclassical implications. We also have obtained direct evidence for the validity of unitary dynamics and the superposition principle in all experiments conducted so far, although this has forced us to again accept extremely nonclassical situations as physical reality. Why not let these experiences guide us to extend our willingness to entirely give up classical prejudice and instead explore the consequences of a strictly unitary quantum theory embedded into a minimal interpretive framework? After all, exploring the implications of pure quantum features to the largest possible extent can in turn give us back the familiar โclassicalโ notions of the world of our experience. As we have discussed in this paper, consequences of highly nonlocal quantum entanglement lead to the local disappearance of quantum interference effects, may explain the origin of probabilities in quantum mechanics, and are likely capable of accounting for the objectification of observables and therefore the emergence of effective classical reality, thus supplying the missing pieces of the basic Everett theory that have frequently been been used to challenge the viability of a relative stateโtype โminimal interpretation.โ
###### Acknowledgements.
The author thanks A. Fine and H. D. Zeh for fruitful discussions and helpful comments.
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# Reduced Decompositions and Permutation Patterns
## 1. Introduction
Reduced decompositions of permutations are classical objects in combinatorics that appear throughout the literature. Following the work of Rodica Simion and Frank Schmidt in , the study of permutation patterns, particularly pattern avoidance, has become a frequently studied field as well.
In , Sara Billey, William Jockusch, and Richard Stanley relate these two concepts, possibly for the first time. There they show that $`321`$-avoiding permutations are exactly those permutations where the subsequence $`i(i\pm 1)i`$ never occurs in a reduced decomposition. Relatedly, Victor Reiner shows in that the number of $`i(i\pm 1)i`$ occurrences in reduced decompositions of the longest element in the symmetric group, which has the maximal number of occurrences of $`321`$, is equal to the number of such reduced decompositions. Stanley had previously shown that this is the number of standard Young tableaux of a staircase shape in .
Inspired by these results, and more generally by the relationship they suggest between the two aspects of permutations, this paper studies elements of the symmetric group from the combined perspectives of their reduced decompositions and their patterns. While these aspects of a permutation appear extensively in combinatorial literature, they are not often treated together. This paper strives to remedy that fact, addressing several questions where reduced decompositions and permutation patterns together lead to interesting results.
After introducing basic terminology and notation in Section 2, Section 3 generalizes the result of Billey, Jockusch, and Stanley, via a new characterization of vexillary permutations in Theorem 3.8. This characterization is based on the reduced decompositions of the permutations *containing* the permutation in question, and is strikingly different from all previous equivalent characterizations. In addition to requiring that each of the permutations containing the vexillary permutation has a certain kind of reduced decomposition, the proof of Theorem 3.8 explicitly constructs such a reduced decomposition.
There are three algorithms which appear in this paper, the first of which occurs in the proof of Theorem 3.8. It should be noted that these are not deterministic, and include a certain amount of choice. For instance, Example 3.9 describes only one possible route that the algorithm VEX may take on a particular input.
There is an equivalence relation, sometimes known as the commutation relation, on the set of reduced decompositions of a particular permutation. This and an associated graph are discussed in Section 4. Theorem 4.6 and Corollary 4.8 characterize permutations with graphs and commutation classes having certain properties. These results are strengthened in Theorem 6.12.
The results in Sections 5 and 6 discuss permutation patterns with respect to a polygon defined by Serge Elnitsky in . The rhombic tilings of this polygon are in bijection with the commutation classes of a permutation. New results include that the number of commutation classes of a permutation is monotonically increasing with respect to pattern containment (Theorem 5.10), and several results pertaining to a poset associated with tilings of the polygon. Finally, Section 7 completely describes this poset in the case of a freely braided permutation, as defined by Richard Green and Jozsef Losonczy in and .
## 2. Basic Definitions
The main definitions and notation that appear throughout the paper are discussed below. For more information about these objects, including proofs of elementary facts, see and .
Let $`๐_n`$ denote the symmetric group on $`n`$ elements. An element $`w๐_n`$ permutes $`\{1,\mathrm{},n\}`$ by mapping $`iw(i)`$. This permutation will be written in one-line notation $`w=w(1)w(2)\mathrm{}w(n)`$.
###### Example 2.1.
$`4213๐_4`$ maps $`1`$ to $`4`$, $`2`$ to itself, $`3`$ to $`1`$, and $`4`$ to $`3`$.
For $`i\{1,\mathrm{},n1\}`$, the map $`s_i`$ transposes $`i`$ and $`i+1`$, and fixes all other elements in a permutation. The symmetric group $`๐_n`$ is the Coxeter group of type $`A_{n1}`$, and it is generated by the adjacent transpositions $`\{s_i:i=1,\mathrm{},n1\}`$. The adjacent transpositions satisfy the Coxeter relations:
(1) $`s_i^2=1`$ $`\text{for all }i;`$
$`s_is_j=s_js_i`$ $`\text{if }|ij|>1;`$
(2) $`s_is_{i+1}s_i=s_{i+1}s_is_{i+1}`$ $`\text{for }1in2.`$
Equation (1) is called the *short braid relation*, and equation (2) is the *long braid relation*. A map is written to the left of its input, so $`s_iw`$ interchanges the positions of values $`i`$ and $`i+1`$ in the permutation $`w`$, while $`ws_i`$ interchanges the values in positions $`i`$ and $`i+1`$ in $`w`$. If $`w=w(1)\mathrm{}w(n)`$, then $`ws_i=w(1)\mathrm{}w(i+1)w(i)\mathrm{}w(n)`$.
Because the symmetric group is generated by adjacent transpositions, any permutation $`w๐_n`$ can be written as $`w=s_{i_1}\mathrm{}s_i_{\mathrm{}}`$ for some $`\{i_1,\mathrm{},i_{\mathrm{}}\}`$. The least such $`\mathrm{}`$ is the *length* of $`w`$, denoted $`\mathrm{}(w)`$. An *inversion* in $`w`$ is a pair $`(i,j)`$ where $`i<j`$ and $`w(i)>w(j)`$. The *inversion set* is $`I(w)=\{(i,j):(i,j)\text{ is an inversion}\}`$. Since $`I(w)[1,n]\times [1,n]`$, the inversion set can also be viewed as an array. The number of inversions in $`w`$ is equal to $`\mathrm{}(w)`$ (see ). For obvious reasons, the permutation $`w_0:=n\mathrm{}21๐_n`$ is called the *longest element* in $`๐_n`$.
###### Definition 2.2.
For a permutation $`w`$ with $`\mathrm{}(w)=\mathrm{}`$, a string $`i_1\mathrm{}i_{\mathrm{}}`$ such that $`w=s_{i_1}\mathrm{}s_i_{\mathrm{}}`$ is a *reduced decomposition* of $`w`$. (Some sources call this a *reduced word*.) The set $`R(w)`$ consists of all reduced decompositions of $`w`$.
###### Definition 2.3.
A *factor* is a consecutive substring of a reduced decomposition.
Similar to the Coxeter relations, a factor $`j_1j_2`$ in a reduced decomposition will be called a *short braid move* if $`|j_1j_2|>1`$, and a factor $`j(j\pm 1)j`$ will be called a *long braid move*. The set $`R(w)`$ has been studied in various contexts, notably by Stanley in . There, Stanley computes $`|R(w)|`$ for several classes of permutations in terms of the number of standard Young tableaux of certain shapes. In the case of a vexillary permutation, this is equal to the number of standard Young tableaux of a single shape $`\lambda (w)`$ (see also Exercise 7.22 of ). The definition of vexillary permutations is postponed until Section 3, where they will be discussed in depth.
###### Definition 2.4.
Let $`w=w(1)\mathrm{}w(n)`$ and $`p=p(1)\mathrm{}p(k)`$ for $`kn`$. The permutation $`w`$ *contains the pattern $`p`$* if there exist $`i_1<\mathrm{}<i_k`$ such that $`w(i_1)\mathrm{}w(i_k)`$ is in the same relative order as $`p(1)\mathrm{}p(k)`$. That is, $`w(i_h)<w(i_j)`$ if and only if $`p(h)<p(j)`$. If $`w`$ does not contain $`p`$, then $`w`$ *avoids* $`p`$, or is *$`p`$-avoiding*.
Suppose that $`w`$ contains the pattern $`p`$, with $`\{i_1,\mathrm{},i_k\}`$ as defined above. Then $`w(i_1)\mathrm{}w(i_k)`$ is an *occurrence* of $`p`$ in $`w`$. The notation $`p(j)`$ will denote the value $`w(i_j)`$. If $`\overline{p}=p(j)p(j+1)\mathrm{}p(j+m)`$, then $`\overline{p}=w(i_j)w(i_{j+1})\mathrm{}w(i_{j+m})`$.
###### Example 2.5.
Let $`w=7413625`$, $`p=1243`$, and $`q=1234`$. Then $`1365`$ is an occurrence of $`p`$, with $`1=1`$, $`2=3`$, $`4=6`$, and $`3=5`$. Also, $`24=36`$. The permutation $`w`$ is $`q`$-avoiding.
###### Definition 2.6.
Let $`w`$ contain the pattern $`p`$, and let $`p`$ be a particular occurrence of $`p`$. If $`w(j)p`$, then $`w(j)`$ is a *pattern entry* in $`w`$. Otherwise $`w(j)`$ is a *non-pattern entry*. If a non-pattern entry lies between two pattern entries in the one-line notation for $`w`$, then it is *inside* the pattern. Otherwise it is *outside* the pattern. โInsideโ and โoutsideโ are only defined for non-pattern entries.
###### Definition 2.7.
Let $`p`$ be an occurrence of $`p๐_k`$ in $`w`$. Suppose that $`x`$ is inside the pattern, that $`m<x<m+1`$ for some $`m[1,k1]`$, and that the values $`\{m,x,m+1\}`$ appear in increasing order in the one-line notation for $`w`$. Let $`a,b`$ be maximal so that the values
$$\{ma,ma+1,\mathrm{},m,x,m+1,\mathrm{},m+b1,m+b\}$$
appear in increasing order in the one-line notation for $`w`$. The entry $`x`$ is *obstructed to the left* if a pattern entry smaller than $`ma`$ appears between $`ma`$ and $`x`$ in $`w`$. Likewise, $`x`$ is *obstructed to the right* if a pattern entry larger than $`m+b`$ appears between $`x`$ and $`m+b`$ in $`w`$.
###### Example 2.8.
Let $`w=32451`$ and $`p=3241`$. Then $`3241`$ and $`3251`$ are both occurrences of $`p`$ in $`w`$. Obstruction is only defined for the latter, with $`x=4`$ and $`m=3`$. Then $`a=b=0`$, and $`4`$ is obstructed to the left and not to the right.
###### Example 2.9.
Let $`w=21354`$ and $`p=2143`$. Then $`2154`$ is an occurrence of $`p`$ in $`w`$. Using $`x=3`$, $`m=2`$ in Definition 2.7 shows that $`a=b=0`$, and $`3`$ is obstructed both to the left and to the right.
## 3. Vexillary Characterization
Vexillary permutations first appeared in and subsequent publications by Alain Lascoux and Marcel-Paul Schรผtzenberger. They were also independently found by Stanley in . There have since emerged several equivalent definitions of these permutations, and a thorough discussion of these occurs in . The original definition of Lascoux and Schรผtzenberger, and the one of most relevance to this discussion, is the following.
###### Definition 3.1.
A permutation is *vexillary* if it is $`2143`$-avoiding.
###### Example 3.2.
The permutation $`3641572`$ is vexillary, but $`3641752`$ is not vexillary because $`3175`$ is an occurrence of $`2143`$ in the latter.
The following proposition is key to proving one direction of Theorem 3.8.
###### Proposition 3.3.
Let $`w`$ contain the pattern $`p`$. Let $`x`$ be inside the pattern, with $`m<x<m+1`$ and the values $`\{m,x,m+1\}`$ appearing in increasing order in $`w`$. If $`p`$ is vexillary then $`x`$ cannot be obstructed both to the left and to the right.
###### Proof.
Such obstructions would create a $`2143`$-pattern in $`p`$. โ
Example 2.9 illustrates a non-vexillary permutation which has an element $`x`$ that is obstructed on both sides.
Equivalent characterizations of vexillarity concern the inversion set $`I(w)`$ or the following objects.
###### Definition 3.4.
The *diagram* of a permutation $`w`$ is $`D(w)[1,n]\times [1,n]`$ where
$$(i,j)D(w)\text{ if and only if }i<w^1(j)\text{ and }j<w(i).$$
###### Definition 3.5.
The *code* of $`w`$ is the vector $`c(w)=(c_1(w),\mathrm{},c_n(w))`$ where $`c_i(w)`$ is the number of elements in row $`i`$ of $`I(w)`$. The shape $`\lambda (w)`$ is the partition formed by writing the entries of the code in non-increasing order.
###### Proposition 3.6.
The following are equivalent definitions of vexillarity for a permutation $`w`$:
1. $`w`$ is $`2143`$-avoiding;
2. The set of rows of $`I(w)`$ is totally ordered by inclusion;
3. The set of columns of $`I(w)`$ is totally ordered by inclusion;
4. The set of rows of $`D(w)`$ is totally ordered by inclusion;
5. The set of columns of $`D(w)`$ is totally ordered by inclusion;
6. $`\lambda (w)^{}=\lambda (w^1)`$, where $`\lambda (w)^{}`$ is the transpose of $`\lambda (w)`$.
###### Proof.
See . โ
This section proves a new characterization of vexillary permutations, quite different from those in Proposition 3.6. A partial ordering can be placed on the set of all permutations $`๐_1๐_2๐_3\mathrm{}`$, where $`u<v`$ if $`v`$ contains the pattern $`u`$. Definition 3.1 determines vexillarity by a condition on the principal order ideal of a permutation. The new characterization, Theorem 3.8, depends on a particular condition holding for the principal *dual* order ideal.
###### Definition 3.7.
Let $`๐=i_1\mathrm{}i_{\mathrm{}}`$ be a reduced decomposition of $`w=w(1)\mathrm{}w(n)`$. For $`M`$, the *shift of $`๐ข`$ by $`M`$* is
$$๐^M:=(i_1+M)\mathrm{}(i_{\mathrm{}}+M)R\left(12\mathrm{}M(w(1)+M)(w(2)+M)\mathrm{}(w(n)+M)\right).$$
###### Theorem 3.8.
The permutation $`p`$ is vexillary if and only if, for every permutation $`w`$ containing a $`p`$-pattern, there exists a reduced decomposition $`๐ฃR(w)`$ containing some shift of an element $`๐ขR(p)`$ as a factor.
###### Proof.
First suppose that $`p๐_k`$ is vexillary. Let $`w๐_n`$ contain a $`p`$-pattern. Assume for the moment that there is a
(3)
$$\stackrel{~}{w}=\left(s_{I_1}\mathrm{}s_{I_q}\right)w\left(s_{J_1}\mathrm{}s_{J_r}\right)๐_n$$
such that
1. $`\mathrm{}(\stackrel{~}{w})=\mathrm{}(w)(q+r)`$;
2. $`\stackrel{~}{w}`$ has a $`p`$-pattern in positions $`\{1+M,\mathrm{},k+M\}`$ for some $`M[0,nk]`$.
Choose a reduced decomposition $`๐R(p)`$. Let $`\stackrel{~}{w}^{}๐_n`$ be the permutation obtained from $`\stackrel{~}{w}`$ by placing the values $`\{\stackrel{~}{w}(1+M),\mathrm{},\stackrel{~}{w}(k+M)\}`$ in increasing order and leaving all other entries unchanged. Choose any $`๐R(\stackrel{~}{w}^{})`$. Then
(4)
$$\left(I_q\mathrm{}I_1\right)๐๐^M\left(J_r\mathrm{}J_1\right)R(w).$$
It remains only to find a $`\stackrel{~}{w}๐_n`$ satisfying (R$`1`$) and (R$`2`$). This will be done by an algorithm VEX that takes as input a permutation $`w๐_n`$ containing a $`p`$-pattern and outputs the desired permutation $`\stackrel{~}{w}๐_n`$. Because the details of this algorithm can be cumbersome, a brief description precedes each of the major steps.
* Algorithm VEX
INPUT: $`w๐_n`$ with an occurrence $`p`$ of the pattern $`p๐_k`$.
OUTPUT: $`\stackrel{~}{w}๐_n`$ as in equation (3) satisfying (R1) and (R2).
1. Initialize variables.
Set $`w_{[0]}:=w`$ and $`i:=0`$.
2. Check if ready to output.
If $`w_{[i]}`$ has no entries inside the pattern, then OUTPUT $`w_{[i]}`$. Otherwise, choose $`x_{[i]}`$ inside the pattern.
3. Move all inside entries larger than $`k`$ to the right of $`p`$.
If $`x_{[i]}>k`$, then BEGIN
1. Let $`B(x_{[i]})=\{yx_{[i]}:y\text{ is inside the pattern}\}`$.
2. Consider the elements of $`B(x_{[i]})`$ in decreasing order. Multiply $`w_{[i]}`$ on the right by adjacent transpositions (changing *positions* in the one-line notation) to move each element immediately to the right of $`p`$.
3. Let $`w_{[i+1]}`$ be the resulting permutation. Set $`i:=i+1`$ and GOTO Step 1.
4. Move all inside entries smaller than $`1`$ to the left of $`p`$.
If $`x_{[i]}<1`$, then BEGIN
1. Let $`S(x_{[i]})=\{yx_{[i]}:y\text{ is inside the pattern}\}`$.
2. Consider the elements of $`S(x_{[i]})`$ in increasing order. Multiply $`w_{[i]}`$ on the right by adjacent transpositions to move each element immediately to the left of $`p`$.
3. Let $`w_{[i+1]}`$ be the resulting permutation. Set $`i:=i+1`$ and GOTO Step 1.
5. Determine bounds in the pattern for the inside entry.
Let $`m[1,k1]`$ be the unique value such that $`m<x_{[i]}<m+1`$.
6. Change the occurrence of $`p`$ and the inside entry so that it does not lie between its bounds in the pattern.
If the values $`\{m,x_{[i]},m+1\}`$ appear in increasing order in $`w_{[i]}`$, then define $`a`$ and $`b`$ as in Definition 2.7 and BEGIN
1. If $`x_{[i]}`$ is unobstructed to the right, then BEGIN
1. Let $`R(x_{[i]})`$ be the set of non-pattern entries at least as large as $`x_{[i]}`$ and lying between $`x_{[i]}`$ and $`m+b`$ in the one-line notation of $`w_{[i]}`$.
2. Consider the elements of $`R(x_{[i]})`$ in decreasing order. For each $`yR(x_{[i]})`$, multiply on the right by adjacent transpositions until $`y`$ is immediately to the right of $`m+b`$, or the right neighbor of $`y`$ is $`z>y`$. In the latter case, $`z=m+b^{}`$ for some $`b^{}[1,b]`$ because all larger non-pattern entries are already to the right of $`m+b`$. Interchange the roles of $`y`$ and $`m+b^{}`$, and move this new $`y`$ to the right in the same manner, until it is to the right of (the redefined) $`m+b`$.
3. Let $`w_{[i+1]}`$ be the resulting permutation, with $`p`$ redefined as indicated. Let $`x_{[i+1]}`$ be the non-pattern entry in the final move after any interchange of roles. This is greater than $`x_{[i]}`$ and the newly redefined $`m+b`$, and occurs to the right of the new $`m+b`$. If $`x_{[i+1]}`$ is outside of the pattern, GOTO Step 1 with $`i:=i+1`$. Otherwise GOTO Step 2 with $`i:=i+1`$.
2. The entry $`x_{[i]}`$ is unobstructed to the left (Proposition 3.3). BEGIN
1. Let $`L(x_{[i]})`$ be the set of non-pattern entries at most as large as $`x_{[i]}`$ and lying between $`ma`$ and $`x_{[i]}`$ in the one-line notation of $`w_{[i]}`$.
2. Consider the elements of $`L(x_{[i]})`$ in increasing order. For each $`yL(x_{[i]})`$, multiply on the right by adjacent transpositions until $`y`$ is immediately to the left of $`ma`$, or the left neighbor of $`y`$ is $`z<y`$. In the latter case, $`z=ma^{}`$ for some $`a^{}[0,a]`$ because all smaller non-pattern entries are already to the left of $`ma`$. Interchange the roles of $`y`$ and $`ma^{}`$, and move this new $`y`$ to the left in the same manner, until it is to the left of (the redefined) $`ma`$.
3. Let $`w_{[i+1]}`$ be the resulting permutation, with $`p`$ redefined as indicated. Let $`x_{[i+1]}`$ be the non-pattern entry in the final move after any interchange of roles. This is less than $`x_{[i]}`$ and the newly redefined $`ma`$, and occurs to the left of the new $`ma`$. If $`x_{[i+1]}`$ is outside of the pattern, GOTO Step 1. Otherwise GOTO Step 3 with $`i:=i+1`$.
7. Change the occurrence of $`p`$, but not its position, so that the value of the inside entry increases but the values of $`p`$ either stay the same or decrease.
If $`w_{[i]}(s)=m+1`$ and $`w_{[i]}(t)=x_{[i]}`$ with $`s<t`$, multiply $`w_{[i]}`$ on the left by adjacent transpositions (changing *values* in the one-line notation) to obtain $`w_{[i+1]}`$ with the values $`[x_{[i]},m+1]`$ in increasing order. Then $`w_{[i+1]}(s)`$ is in the half-open interval $`[x_{[i]},m+1)`$, and $`w_{[i+1]}(t)`$ is in the half-open interval $`(x_{[i]},m+1]`$. GOTO Step 2 with $`x_{[i+1]}:=w_{[i+1]}(t)`$, the pattern redefined so that $`m+1:=w_{[i+1]}(s)`$, and $`i:=i+1`$.
8. Change the occurrence of $`p`$, but not its position, so that the value of the inside entry decreases but the values of $`p`$ either stay the same or increase.
If $`w_{[i]}(s)=m`$ and $`w_{[i]}(t)=x_{[i]}`$ with $`s>t`$, multiply $`w_{[i]}`$ on the left by adjacent transpositions to obtain $`w_{[i+1]}`$ with the values $`[m,x_{[i]}]`$ in increasing order. Then $`w_{[i+1]}(s)`$ is in the half-open interval $`(m,x_{[i]}]`$, and $`w_{[i+1]}(t)`$ is in the half-open interval $`[m,x_{[i]})`$. GOTO Step 3 with $`x_{[i+1]}:=w_{[i+1]}(t)`$, the pattern redefined so that $`m:=w_{[i+1]}(s)`$, and $`i:=i+1`$.
Each subsequent visit to Step 1 involves a permutation with strictly fewer entries inside the pattern than on the previous visit. Each multiplication by an adjacent transposition indicated in the algorithm removes an inversion, and so decreases the length of the permutation. This is crucial because of requirement (R$`1`$).
Consider the progression of VEX:
* Step 1 $``$ HALT or begin a pass through VEX;
* Step 2 $``$ Step 1;
* Step 3 $``$ Step 1;
* Step 5a $``$ Steps 2 or 6;
* Step 5b $``$ Steps 3 or 7;
* Step 6 $``$ Steps 2, 5, or 6;
* Step 7 $``$ Steps 3, 5, or 7.
Step 5a concludes with $`x_{[i+1]}`$ to the left of its lower pattern bound, and smaller pattern elements lying between $`x_{[i+1]}`$ and this bound. Therefore, no matter how often Step 6 is next called, the algorithm will never subsequently go to Step 5b before going to Step 1. Likewise, a visit to Step 5b means that Step 5a can never be visited until Step 1 is visited and a new entry inside the pattern is chosen.
Steps 2 and 3 do not change the relative positions of $`p`$.
Steps 5a and 6 imply $`x_{[i+1]}>x_{[i]}`$, while $`x_{[i+1]}<x_{[i]}`$ after Steps 5b and 7. Let $`m`$ be as in Step 4. Until revisiting Step 1, the values $`m^{}`$, for $`m^{}m+1`$, do not increase if $`x_{[i+1]}>x_{[i]}`$. Nor do the values $`m^{}`$, for $`m^{}m`$, decrease if $`x_{[i+1]}<x_{[i]}`$. The other pattern values are unchanged. The definition of $`m`$ means that the reordering of values in Steps 6 and 7 does not change the positions in which the pattern $`p`$ occurs. Additionally, these steps change the value of the entry inside the pattern (that is, $`x_{[i+1]}x_{[i]}`$), but not its position.
These observations indicate not only that VEX terminates, but that it outputs $`\stackrel{~}{w}๐_n`$ as in equation (3) satisfying (R1) and (R2). This completes one direction of the proof.
Now suppose $`p๐_k`$ is not vexillary. There is an occurrence $`2143`$ such that
$$p=\mathrm{}2\mathrm{}1(2+1)(2+2)\mathrm{}(32)(31)4\mathrm{}3\mathrm{}.$$
Define $`z`$ to be the index such that $`p(z)=1`$. Define $`w๐_{k+1}`$ by
$$w(m)=\{\begin{array}{cc}\hfill p(m):& mz\text{ and }p(m)2;\hfill \\ \hfill p(m)+1:& mz\text{ and }p(m)>2;\hfill \\ \hfill 2+1:& m=z+1;\hfill \\ \hfill p(m1):& m>z+1\text{ and }p(m)2;\hfill \\ \hfill p(m1)+1:& m>z+1\text{ and }p(m)>2.\hfill \end{array}$$
For example, if $`p=2143`$, then $`w=21354`$.
If there is a reduced decomposition $`๐R(w)`$ such that $`๐=๐_\mathrm{๐}๐^M๐_\mathrm{๐}`$ for $`๐R(p)`$ and $`M`$, then there is a $`\stackrel{~}{w}๐_{k+1}`$ as in equation (3) satisfying (R1) and (R2). Keeping the values $`1`$, $`2`$, $`3`$, and $`4`$ as defined above, the permutation $`w`$ was constructed so that
$$w=\mathrm{}2\mathrm{}1(2+1)(2+2)\mathrm{}(32)(31)3(4+1)\mathrm{}(3+1)\mathrm{}.$$
One of the values in the consecutive subsequence $`(2+1)(2+2)\mathrm{}3`$ must move to get a consecutive $`p`$-pattern in $`\stackrel{~}{w}`$. However, the values $`\{2,\mathrm{},3+1\}`$ appear in increasing order in $`w`$, and the consecutive subsequence
$$1(2+1)(2+2)\mathrm{}(32)(31)3(4+1)$$
in $`w`$ is increasing. Therefore, there is no way to multiply $`w`$ by adjacent transpositions, always eliminating an inversion, to obtain a consecutive $`p`$-pattern.
Hence, if $`p`$ is not vexillary then there exists a permutation $`w`$ containing a $`p`$-pattern such that no reduced decomposition of $`w`$ contains a shift of a reduced decomposition of $`p`$ as a factor. โ
###### Example 3.9.
If $`w=\mathrm{๐}14\mathrm{๐}5\mathrm{๐}`$ and $`p=231`$, with the chosen occurrence $`p`$ in bold, the algorithm VEX may proceed as follows.
* $`w_{[0]}:=\mathrm{๐}14\mathrm{๐}5\mathrm{๐}`$.
* Step 1: $`x_{[0]}:=1`$.
* Step 3: $`w_{[0]}w_{[0]}s_1=1\mathrm{๐}4\mathrm{๐}5\mathrm{๐}=:w_{[1]}`$.
* Step 1: $`x_{[1]}:=5`$.
* Step 6: $`w_{[1]}s_5w_{[1]}=1\mathrm{๐}4\mathrm{๐}6\mathrm{๐}=:w_{[2]}`$; $`x_{[2]}:=6`$.
* Step 2: $`w_{[2]}w_{[2]}s_5=1\mathrm{๐}4\mathrm{๐๐}6=:w_{[3]}`$.
* Step 1: $`x_{[3]}:=4`$.
* Step 5a: $`w_{[3]}w_{[3]}=1\mathrm{๐๐}5\mathrm{๐}6=:w_{[4]}`$; $`x_{[4]}:=5`$.
* Step 2: $`w_{[4]}w_{[4]}s_4=1\mathrm{๐๐๐}56=:w_{[5]}`$.
* Step 1: output $`1\mathrm{๐๐๐}56`$.
Therefore $`\stackrel{~}{w}=134256=s_5ws_1s_5s_4`$, and $`\stackrel{~}{w}^{}=123456`$. Keeping the notation of equation (4), $`๐=\mathrm{}`$ and $`M=1`$. The unique reduced decomposition of $`231`$ is $`12`$, and indeed
$$(5)\mathrm{}(12)^1(451)=523451R(w).$$
###### Example 3.10.
Let $`w=21354`$ and $`p=2143`$. No element of $`R(w)=\{14,41\}`$ contains a shift of any element of $`R(p)=\{13,31\}`$ as a factor.
###### Remark 3.11.
Suppose that $`๐R(w)`$ contains a shift of $`๐R(p)`$ as a factor,
(5)
$$๐=๐_\mathrm{๐}๐^M๐_\mathrm{๐}.$$
Then $`๐R(p)`$ can be replaced by any $`๐^{\mathbf{}}R(p)`$ in equation (5).
Some care must be taken regarding factors in reduced decompositions. This is clarified in the following definition and lemma, the proof of which is straightforward.
###### Definition 3.12.
Let $`w๐_n`$ and $`๐R(w)`$. Write $`๐=๐๐๐`$, where $`๐R(u)`$ and $`๐R(v)`$. Suppose that $`๐`$ contains only letters in $`S=\{1+M,\mathrm{},k1+M\}`$. If no element of $`R(u)`$ has an element of $`S`$ as its rightmost character and no element of $`R(v)`$ has an element of $`S`$ as its leftmost character, then $`๐`$ is *isolated* in $`๐`$. Equivalently, the values $`\{1+M,\mathrm{},k+M\}`$ must appear in increasing order in $`v`$, and the positions $`\{1+M,\mathrm{},k+M\}`$ must comprise an increasing sequence in $`u`$.
If $`๐R(k\mathrm{}21)`$ and a shift of $`๐`$ appears as a factor in a reduced decomposition of some permutation, then $`๐`$ is necessarily isolated. This is because $`๐`$ has maximal reduced length in the letters $`\{1,\mathrm{},k1\}`$, so any factor of length greater than $`\left(\genfrac{}{}{0pt}{}{k}{2}\right)`$ in the letters $`\{1+M,\mathrm{},k1+M\}`$ is not reduced.
###### Lemma 3.13.
If a reduced decomposition of $`w`$ contains an isolated shift of a reduced decomposition of $`p`$, then $`w`$ contains the pattern $`p`$.
The converse to Lemma 3.13 holds if $`p`$ is vexillary.
The characterization of vexillary in Theorem 3.8 differs substantially from those in Proposition 3.6. There is not an obvious way to prove equivalence with any of the definitions (V2)-(V6), except via (V1). This raises the question of whether more may be understood about vexillary permutations (or perhaps other types, such as Grassmannian or dominant permutations) by studying their reduced decompositions or the permutations that contain those in question as patterns.
Theorem 3.8 has a number of consequences, and will be used often in the subsequent sections of this paper. Most immediately, notice that it generalizes the result of Billey, Jockusch, and Stanley mentioned earlier: $`321`$-avoiding permutations are exactly those whose reduced decompositions contain no long braid moves, and observe that $`R(321)=\{121,212\}`$.
## 4. The Commutation Relation
Recall the definition of short and long braid moves in a reduced decomposition, as well as the short and long braid relations described in equations (1) and (2). It is well known that any element of $`R(w)`$ can be transformed into any other element of $`R(w)`$ by successive applications of the braid relations.
Because the short braid relation represents the commutativity of particular pairs of adjacent transpositions, the following equivalence relation is known as the *commutation relation*.
###### Definition 4.1.
For a permutation $`w`$ and $`๐,๐R(w)`$, write $`๐๐`$ if $`๐`$ can be obtained from $`๐`$ by a sequence of short braid moves. Let $`C(w)`$ be the set of commutation classes of reduced decompositions of $`w`$, as defined by $``$.
###### Example 4.2.
The commutation classes of $`4231๐_4`$ are $`\{12321\}`$, $`\{32123\}`$, and $`\{13231,31231,13213,31213\}`$.
###### Definition 4.3.
For a permutation $`w`$, the graph $`G(w)`$ has vertex set equal to $`C(w)`$, and two vertices share an edge if there exist representatives of the two classes that differ by a long braid move.
Elnitsky gives a very elegant representation of this graph in , which will be discussed in depth in Section 5. A consequence of his description, although not difficult to prove independent of his work, is the following.
###### Proposition 4.4.
The graph $`G(w)`$ is connected and bipartite.
###### Proof.
See . โ
Despite Proposition 4.4, much remains to be understood about the graph $`G(w)`$. For example, even the size of the graph for $`w_0`$ (that is, the number of commutation classes for the longest element) is unknown.
Billey, Jockusch, and Stanley characterize all permutations with a single commutation class, and hence whose graphs are a single vertex, as $`321`$-avoiding permutations. A logical question to ask next is: for what permutations does each reduced decomposition contain at most one long braid move? More restrictively: what if this long braid move is required to be a specific shift of $`121`$ or $`212`$? Moreover, what are the graphs in these cases?
###### Definition 4.5.
Let $`U_n=\{w๐_n:\text{no }๐R(w)\text{ has two long braid moves}\}`$.
###### Theorem 4.6.
$`U_n`$ is the set of permutations such that every $`321`$-pattern in $`w`$ has the same maximal element and the same minimal element.
###### Proof.
Assume $`w`$ has a $`321`$-pattern. Suppose that every occurrence of $`321`$ in $`w`$ has $`3=x`$ and $`1=y`$. Suppose that $`๐R(w)`$ has at least one long braid move. Choose $`k`$ so that $`j_kj_{k+1}j_{k+2}`$ is the first such. Each adjacent transposition in a reduced decomposition increases the length of the product. Then by the supposition,
$$s_{j_{k+2}}s_{j_{k+1}}s_{j_k}\mathrm{}s_{j_1}w$$
is $`321`$-avoiding, so $`j_{k+3}\mathrm{}j_{\mathrm{}}`$ has no long braid moves. It remains only to consider when $`j_{k+2}j_{k+3}j_{k+4}`$ is also a long braid move. The only possible reduced configurations for such a factor $`j_kj_{k+1}j_{k+2}j_{k+3}j_{k+4}`$ are shifts of $`21232`$ and $`23212`$. If either of these is not isolated, then it is part of a shift of $`212321`$, $`321232`$, $`123212`$, or $`232123`$. Notice that
* $`212321,321232,123212,232123R(4321)`$;
* $`23212R(4312)`$;
* $`21232R(3421)`$.
If $`j_kj_{k+1}j_{k+2}j_{k+3}j_{k+4}`$ is isolated in $`๐`$ then $`w`$ contains a $`4312`$\- or $`3421`$-pattern by Lemma 3.13. Otherwise, $`w`$ contains a $`4321`$-pattern. However, every $`321`$-pattern in $`w`$ has $`3=x`$ and $`1=y`$. Therefore $`j_kj_{k+1}j_{k+2}`$ is the only long braid move in $`๐`$, so $`wU_n`$.
Now let $`w`$ be an element of $`U_n`$. If $`w`$ has two $`321`$-patterns that do not have the same maximal element and the same minimal element, then they intersect at most once or they create a $`4321`$-, $`4312`$-, or $`3421`$-pattern. These three patterns are vexillary. Thus by Theorem 3.8 and the examples above, containing one of these patterns would imply that some element of $`R(w)`$ has more than one long braid move. If the two $`321`$-patterns intersect at most once, their union may be a non-vexillary pattern, so Theorem 3.8 does not necessarily apply. However, a case analysis shows that it is possible to shorten $`w`$ by adjacent transpositions and make one $`321`$-pattern increasing (via a long braid move) without destroying the other $`321`$-pattern. Thus an element of $`R(w)`$ would have more than one long braid move, contradicting $`wU_n`$. โ
###### Definition 4.7.
Let $`U_n(j)`$ consist of permutations with some $`321`$-pattern, where every long braid move that occurs must be $`j(j+1)j`$ or $`(j+1)j(j+1)`$.
###### Corollary 4.8.
$`U_n(j)=\{w๐_n:w`$ has a unique $`321`$-pattern and $`2=j+1\}`$. If $`w`$ has a unique $`321`$-pattern, then $`w(2)=2`$.
###### Proof.
A unique $`321`$-pattern implies that $`\{1,\mathrm{},21\}1`$ all appear to the left of $`2`$ in $`w(1)\mathrm{}w(n)`$, and $`\{2+1,\mathrm{},n\}3`$ all appear to the right of $`2`$, so the second statement follows.
Consider the long braid moves that may appear for elements of $`U_nU_n(j)`$. Let $`wU_n`$ have $`k`$ distinct $`321`$-patterns. By Theorem 4.6, these form a pattern $`p=(k+2)23\mathrm{}k(k+1)1๐_{k+2}`$ in $`w`$. The permutation $`p`$ is vexillary, so there exists $`M`$ and a reduced decomposition $`๐_\mathrm{๐}๐^M๐_\mathrm{๐}R(w)`$ for each $`๐R(p)`$. There are elements in $`R(p)`$ with long braid moves $`i(i+1)i`$ for each $`i[1,k]`$. For example, $`12\mathrm{}k(k+1)k\mathrm{}21R(p)`$. Therefore, if $`wU_n(j)`$, then $`k=1`$, so $`w`$ has a unique $`321`$-pattern.
Suppose that $`w`$ has a unique $`321`$-pattern. Because $`w(2)=2`$, the only possible long braid moves in reduced decompositions of $`w`$ are $`(21)2(21)`$ or $`2(21)2`$. โ
###### Corollary 4.9.
If $`wU_n`$ and $`w`$ has $`k`$ distinct $`321`$-patterns, then $`|C(w)|=k+1`$ and the graph $`G(w)`$ is a path of $`k+1`$ vertices connected by $`k`$ edges.
###### Proof.
Because $`w`$ contains the pattern $`p=(k+2)23\mathrm{}k(k+1)1๐_{k+2}`$, there is a subgraph of $`G(w)`$ that is a path of $`k+1`$ vertices connected by $`k`$ edges. Since $`p`$ accounts for all of the $`321`$-patterns in $`w`$, this is all of $`G(w)`$. โ
###### Corollary 4.10.
If $`wU_n(j)`$, then $`|C(w)|=2`$ and the graph $`G(w)`$ is a pair of vertices connected by an edge.
## 5. Elnitskyโs Polygon
In his doctoral thesis and in , Elnitsky developed a bijection between commutation classes of reduced decompositions of $`w๐_n`$ and rhombic tilings of a particular $`2n`$-gon $`X(w)`$. This bijection leads to a number of interesting questions about tilings of $`X(w)`$ and their relations to the permutation $`w`$ itself. A number of these ideas are studied in this and the following section.
###### Definition 5.1.
For $`w๐_n`$, let $`X(w)`$ be the $`2n`$-gon with all sides of unit length such that
1. Sides of $`X(w)`$ are labeled $`1,\mathrm{},n,w(n),\mathrm{},w(1)`$ in order;
2. The portion labeled $`1,\mathrm{},n`$ is convex; and
3. Sides with the same label are parallel.
Orient the polygon so that the edge labeled $`1`$ lies to the left of the top vertex and the edge labeled $`w(1)`$ lies to its right. This is *Elnitskyโs polygon*.
###### Example 5.2.
For $`w_0๐_n`$, the polygon $`X(w_0)`$ is a centrally symmetric $`2n`$-gon.
###### Definition 5.3.
The hexagon $`X(321)`$ can be tiled by rhombi with sides of unit length in exactly two ways. Each of these tilings is the *flip* of the other.
###### Definition 5.4.
Let $`T(w)`$ be the set of tilings of $`X(w)`$ by rhombi with sides of unit length. Define a graph $`G^{}(w)`$ with vertex set $`T(w)`$, and connect two tilings by an edge if they differ by a flip of the tiling of a single sub-hexagon.
Unless otherwise indicated, the term *tiling* refers to an element of $`T(w)`$.
###### Theorem 5.5 (Elnitsky).
The graphs $`G(w)`$ and $`G^{}(w)`$ are isomorphic.
Henceforth, both graphs will be denoted $`G(w)`$.
Before discussing new results related to this polygon, it is important to understand Elnitskyโs bijection, outlined in the following algorithm. A more thorough treatment appears in .
* Algorithm ELN
INPUT: $`TT(w)`$.
OUTPUT: An element of $`C_TC(w)`$.
1. Set the polygon $`P_{[0]}:=X(w)`$, the string $`๐_{\mathbf{[}\mathrm{๐}\mathbf{]}}:=\mathrm{}`$, and $`i:=0`$.
2. If $`P_{[i]}`$ has no area, then OUTPUT $`๐_{\mathbf{[}๐\mathbf{]}}`$.
3. There is at least one tile $`t_i`$ that shares two edges with the right side of $`P_{[i]}`$.
4. If $`t_i`$ includes the $`j^{\text{th}}`$ and $`(j+1)^{\text{st}}`$ edges from the top along the right side of $`P_{[i]}`$, set $`๐_{\mathbf{[}๐\mathbf{+}\mathrm{๐}\mathbf{]}}:=j๐_{\mathbf{[}๐\mathbf{]}}`$.
5. Let $`P_{[i+1]}`$ be $`P_{[i]}`$ with the tile $`t_i`$ removed. Set $`i:=i+1`$ and GOTO Step 1.
ELN yields the entire commutation class because of the choice of tile in Step 2.
###### Example 5.6.
The tiling in Figure 2 corresponds to the equivalence class consisting solely of the reduced decomposition $`12343212R(53241)`$.
###### Corollary 5.7.
If $`p`$ is vexillary and $`w`$ contains a $`p`$-pattern, then $`G(p)`$ is a subgraph of $`G(w)`$.
###### Proof.
This follows from Theorems 3.8 and 5.5. โ
Elnitskyโs correspondence, described in ELN, combined with Theorem 4.6 and Corollary 4.8, indicates that any tiling of $`X(w)`$ for $`wU_n`$ has at most one sub-hexagon (every tiling has exactly one sub-hexagon if $`w`$ is not $`321`$-avoiding). Moreover, the sub-hexagon has the same vertical position for all elements of $`U_n(j)`$.
Under certain circumstances, the polygon $`X(w)`$ for $`w๐_n`$ can be rotated or reflected to give a polygon $`X(w^{})`$ for another $`w^{}๐_n`$.
###### Corollary 5.8.
Let $`w=w(1)\mathrm{}w(n)`$ and $`w^R=w(n)\mathrm{}w(1)`$. Then $`|C(w)|=|C(w^R)|`$ and $`G(w)G(w^R)`$.
###### Corollary 5.9.
Let $`w=w(1)\mathrm{}w(n)`$. If $`w(1)=n,w(2)=n1,\mathrm{},w(i)=n+1i`$, then $`|C(w)|=|C(w^{(i)})|`$and $`G(w)G(w^{(i)})`$ where
$$w^{(i)}=(w(i+1)+i)(w(i+2)+i)\mathrm{}(w(n)+i)i(i1)\mathrm{}21$$
and all entries are modulo $`n`$. Likewise, if $`w(n)=1,w(n1)=2,\mathrm{},w(nj+1)=j`$, then $`|C(w)|=|C(w_{(j)})|`$ and $`G(w)G(w_{(j)})`$ where
$$w_{(j)}=n(n1)\mathrm{}(nj+1)(w(1)j)(w(2)j)\mathrm{}(w(nj)j)$$
and all entries are modulo $`n`$.
Elnitskyโs result interprets the commutation classes of $`R(w)`$ as rhombic tilings of $`X(w)`$, with long braid moves represented by flipping sub-hexagons. The following theorem utilizes this interpretation, and demonstrates that the number of commutation classes of a permutation is monotonically increasing with respect to pattern containment, thus generalizing one aspect of Corollary 5.7. Note that $`p`$ is not required to be vexillary in Theorem 5.10, unlike in Theorem 3.8.
###### Theorem 5.10.
If $`w`$ contains the pattern $`p`$, then $`|C(w)||C(p)|`$.
###### Proof.
Consider a tiling $`TT(p)`$. This represents a commutation class of $`R(p)`$. For an ordering of the tiles in $`T`$ as defined by ELN, label the tile $`t_0`$ by $`\mathrm{}(p)`$, the tile $`t_1`$ by $`\mathrm{}(p)1`$, and so on. If the tile with label $`r`$ corresponds to the adjacent transposition $`s_{i_r}`$, then $`i_1\mathrm{}i_{\mathrm{}(p)}R(p)`$.
* Algorithm MONO
INPUT: $`w`$ containing the pattern $`p`$ and $`TT(p)`$ with tiles labeled as described.
OUTPUT: $`T^{}T(w)`$.
1. Set $`w_{[0]}:=w`$, $`p_{[0]}:=p`$, $`T_{[0]}:=T`$, $`T_{[0]}^{}:=\mathrm{}`$, and $`i:=0`$.
2. If $`p_{[i]}`$ is the identity permutation, then define $`T_{[i+1]}^{}`$ to be the tiles of $`T_{[i]}^{}`$ together with any tiling of $`X(w_{[i]})`$. OUTPUT $`T_{[i+1]}^{}`$.
3. Let $`j_{[i]}`$ be such that the tile labeled $`\mathrm{}(p)i`$ includes edges $`p_{[i]}(j_{[i]})`$ and $`p_{[i]}(j_{[i]}+1)`$. Note that $`p_{[i]}(j_{[i]})>p_{[i]}(j_{[i]}+1)`$.
4. Define $`r<s`$ so that $`w_{[i]}(r)=p_{[i]}(j_{[i]})`$ and $`w_{[i]}(s)=p_{[i]}(j_{[i]}+1)`$. Note that $`w_{[i]}(t)`$ is a non-pattern entry for $`t(r,s)`$.
5. Let $`v_{[i]}`$ be the permutation defined by
$$w_{[i]}(t)\{\begin{array}{cc}\hfill w_{[i]}(t):& t<r\text{ or }t>s\hfill \\ \hfill \stackrel{~}{w}_{[i]}(t):& rts\hfill \end{array}$$
where $`(\stackrel{~}{w}_{[i]}(r),\mathrm{},\stackrel{~}{w}_{[i]}(s))`$ is $`\{w_{[i]}(r),\mathrm{},w_{[i]}(s)\}`$ in increasing order.
6. Set $`w_{[i+1]}:=w_{[i]}v_{[i]}`$, and notice that $`\mathrm{}(w_{[i+1]})=\mathrm{}(w_{[i]})\mathrm{}(v_{[i]})`$.
7. The right boundaries of $`X(w_{[i+1]})`$ and $`X(w_{[i]})`$ differ only in the $`r^{\text{th}},\mathrm{},s^{\text{th}}`$ edges, and the left side of this difference (part of the boundary of $`X(w_{[i+1]})`$) is convex. Therefore, this difference has a rhombic tiling $`t_{[i]}`$. Define $`T_{[i+1]}^{}`$ to be the tiles in $`T_{[i]}^{}`$ together with the tiles in $`t_{[i]}`$.
8. Set $`i:=i+1`$ and GOTO Step 1.
The algorithm MONO takes a tiling $`TT(p)`$ and outputs one of possibly several tilings $`T^{}T(w)`$ due to the choice in Steps 1 and 6. A tiling $`T^{}T(w)`$ so obtained can only come from this $`T`$, although possibly with more than one labeling of the tiles. However, this labeling of the tiles merely reflects the choice of a representative from the commutation class, so indeed $`|T(w)||T(p)|`$, and $`|C(w)||C(p)|`$. โ
###### Example 5.11.
Let $`p=31542`$ and $`w=4617352`$. The pattern $`p`$ occurs in $`w`$ as $`p=41752`$. Figure 3 depicts the output of MONO, given the two tilings of $`X(p)`$.
## 6. The Poset of Tilings
Elnitskyโs bijection considers the rhombic tilings of the polygon $`X(w)`$. Rhombi are a special case of a more general class of objects known as zonotopes.
###### Definition 6.1.
A polytope is a *$`d`$-zonotope* if it is the projection of a regular $`n`$-cube onto a $`d`$-dimensional subspace.
Centrally symmetric convex polygons are exactly the $`2`$-zonotopes. These necessarily have an even number of sides.
###### Definition 6.2.
A *zonotopal tiling* of a polygon is a tiling by centrally symmetric convex polygons.
###### Definition 6.3.
Let $`Z(w)`$ be the set of zonotopal tilings of Elnitskyโs polygon. Rhombi are centrally symmetric, so $`T(w)Z(w)`$.
###### Theorem 6.4.
There is a tiling in $`Z(w)`$ containing a $`2k`$-gon with sides parallel to the sides labeled $`i_1<\mathrm{}<i_k`$ if and only if $`i_k\mathrm{}i_1`$ is an occurrence of $`k\mathrm{}1`$ in $`w`$.
###### Proof.
Since the tiles are convex, a $`2k`$-gon in the tiling with sides as described has right side labeled $`i_k,\mathrm{},i_1`$ from top to bottom and left side labeled $`i_1,\mathrm{},i_k`$ from top to bottom. Therefore Elnitskyโs bijection shows that this tile (or rather, any decomposition of it into rhombi) transforms the sequence $`(i_1,\mathrm{},i_k)`$ into $`(i_k,\mathrm{},i_1)`$. Reduced decompositions have minimal length, so no inversions can be โundoneโ by subsequent adjacent transpositions. Therefore $`i_k\mathrm{}i_1`$ must be an occurrence of $`k\mathrm{}1`$ in $`w`$.
Conversely, suppose that $`i_k\mathrm{}i_1`$ is an occurrence of the vexillary pattern $`k\mathrm{}1`$ in $`w`$. For a decreasing pattern, the algorithm VEX can be modified slightly to produce $`\stackrel{~}{w}`$ as in equation (3), where the consecutive occurrence $`k\mathrm{}1`$ is $`i_k\mathrm{}i_1`$. Let $`๐R(k\mathrm{}1)`$ and $`(I_q\mathrm{}I_1)๐๐^M(J_r\mathrm{}J_1)R(w)`$ for $`๐R(\stackrel{~}{w}^{})`$. Removing the rhombi that correspond to $`s_{J_r}\mathrm{}s_{J_1}`$ yields the polygon $`X(\stackrel{~}{w})`$, and the rhombi that correspond to $`๐^M`$ form a sub-$`2k`$-gon with sides parallel to the sides labeled $`\{i_1,\mathrm{},i_k\}`$ in $`X(w)`$. โ
Less specifically, Theorem 6.4 states that a tiling in $`Z(w)`$ can contain a $`2k`$-gon if and only if $`w`$ has a decreasing subsequence of length $`k`$.
Using a group theoretic argument, Pasechnik and Shapiro showed in that no element of $`Z(n\mathrm{}21)`$ consists entirely of hexagons for $`n>3`$. Their result states that at least one rhombus must be present in a hexagonal/rhombic tiling. Kelly and Rottenberg had previously obtained a better bound in terms of arrangements of pseudolines in .
Working with reduced decompositions and Elnitskyโs polygons yields a different proof that no element of $`Z(n\mathrm{}21)`$ can consist of entirely hexagonal tiles for $`n>3`$, and generalizes the result to other types of tiles. Theorem 6.5 is, in a sense, a counterpart to Theorem 6.4.
###### Theorem 6.5.
Let $`w_0`$ be the longest element in $`๐_n`$. There is a tiling $`ZZ(w_0)`$ consisting entirely of $`2k`$-gons if and only if one of the following is true:
1. $`k=2`$; or
2. $`k=n`$.
###### Proof.
If $`k=2`$, the result holds for all $`n`$: every $`X(w_0)`$ can be tiled by rhombi with unit side length. For the remainder of the proof, assume that $`k3`$.
Suppose that there exists $`ZZ(w_0)`$ consisting entirely of $`2k`$-gons. Then there exists $`๐^{M_1}\mathrm{}๐^{M_r}R(w_0)`$, for some $`๐R(k\mathrm{}21)`$. Because $`k\mathrm{}21`$ is the longest element in $`๐_k`$, the factor $`๐^{M_i}`$ can have any of $`\{1+M_i,\mathrm{},k1+M_i\}`$ at either end. Thus, if $`M_i=M_j=M`$ for $`i<j`$, then $`M\pm (k1)\{M_{i+1},\mathrm{},M_{j1}\}`$. Because $`M_i[0,nk]`$, there is at most one shift by $`0`$ and at most one shift by $`nk`$. In fact, there is exactly one of each of these shifts, since all of $`X(w_0)`$ must be tiled and the shifts correspond to the vertical placement of the $`2k`$-gons.
Consider the tile corresponding to $`๐^0`$. This $`2k`$-gon is as high vertically as possible and it is the only tile placed so high. Thus it shares the top vertex and its incident sides with $`X(w_0)`$. That is, two of the tileโs sides are labeled $`1`$ and $`n`$. Similarly, the tile for $`๐^{nk}`$ also has two sides labeled $`1`$ and $`n`$.
By Theorem 6.4, each of these tiles corresponds to a $`k\mathrm{}21`$-pattern with $`k=n`$ and $`1=1`$. However, once the value $`1`$ is to the right of the value $`n`$, it remains to the right as adjacent positions are transposed to lengthen the permutation. Therefore the two patterns, the tiles, and their shifts must be equal: $`0=nk`$.
Indeed, there is always a tiling $`ZZ(w_0)`$ consisting of a single $`2n`$-gon. โ
There is a poset $`P(w)`$ that arises naturally when studying $`Z(w)`$.
###### Definition 6.6.
For a permutation $`w`$, let the poset $`P(w)`$ have elements equal to the zonotopal tilings $`Z(w)`$, partially ordered by reverse edge inclusion.
###### Example 6.7.
In the poset $`P(53241)`$, the tiling in Figure 2 is less than the tiling in Figure 4.
###### Remark 6.8.
For the longest element $`w_0๐_n`$, the poset $`P(w_0)`$ has a maximal element equal to the tiling in $`Z(w_0)`$ that consists of a single $`2n`$-gon.
###### Remark 6.9.
The minimal elements of $`P(w)`$ are the rhombic tilings, which are the vertices of the graph $`G(w)`$. Moreover, edges in the graph $`G(w)`$ correspond to flipping a single sub-hexagon in the tiling. Therefore these edges correspond to the elements of $`P(w)`$ that cover the minimal elements.
The relationships in Remark 6.9 are immediately apparent. Another relationship is not as obvious. This follows from a result of Boris Shapiro, Michael Shapiro, and Alek Vainshtein in .
###### Lemma 6.10 (Shapiro-Shapiro-Vainshtein).
The set of all $`4`$\- and $`8`$-cycles in $`G(w)`$ form a system of generators for the first homology group $`H_1(G(w),/2)`$.
Additionally, Anders Bjรถrner noted that gluing $`2`$-cells into those $`4`$\- and $`8`$-cycles yields a simply connected complex ().
In , Lemma 6.10 is stated only for $`w=w_0`$. However, the proof easily generalizes to all $`w๐_n`$. A straightforward argument demonstrates that a $`4`$-cycle in $`G(w)`$ corresponds to $`ZZ(w)`$ with rhombi and two hexagons, and an $`8`$-cycle corresponds to $`ZZ(w)`$ with rhombi and an octagon. These are exactly the elements of $`P(w)`$ which cover those that correspond to edges of $`G(w)`$.
###### Corollary 6.11.
The elements of $`P(w)`$ that cover the elements (corresponding to edges of $`G(w)`$) covering the minimal elements (corresponding to vertices of $`G(w)`$) correspond to a system of generators for the first homology group $`H_1(G(w),/2)`$.
Little is known about the structure of the graph $`G(w)`$ for arbitrary $`w`$. However, in some cases a description can be given via Theorem 6.4 and Lemma 6.10.
###### Theorem 6.12.
The following statements are equivalent for a permutation $`w`$:
1. $`G(w)`$ is a tree;
2. $`G(w)`$ is a path (that is, no vertex has more than two incident edges);
3. The maximal elements of $`P(w)`$ cover the minimal elements;
4. $`w`$ is $`4321`$-avoiding and any two $`321`$-patterns intersect at least twice.
###### Proof.
(1) $``$ (3) by Corollary 6.11. From Theorem 6.4 and the discussion preceding Corollary 6.11, an $`8`$-cycle in the graph is equivalent to having a $`4321`$-pattern. Similarly, a $`4`$-cycle is equivalent to two sub-hexagons whose intersection has zero area, so some reduced decomposition has two disjoint long braid moves. This implies that two $`321`$-patterns intersect in at most one position. Therefore (1) $``$ (4).
Finally, suppose that $`G(w)`$ is a tree and a vertex has three incident edges. The corresponding tiling has at least three sub-hexagons. However, it is impossible for every pair of these to overlap. This contradicts (1) $``$ (3) $``$ (4), so (1) $``$ (2). โ
If $`C_n`$ is the set of all $`w๐_n`$ for which $`G(w)`$ is a path, then $`U_n(j)U_nC_n`$ by Corollaries 4.9 and 4.10.
Following convention, the unique maximal element in a poset, if it exists, is denoted $`\widehat{1}`$. Remark 6.8 noted that the poset $`P(w_0)`$ has a $`\widehat{1}`$. In fact, there are other $`w`$ for which $`P(w)`$ has a $`\widehat{1}`$, as described below.
###### Theorem 6.13.
The poset $`P(w)`$ has a $`\widehat{1}`$ if and only if $`w`$ is $`4231`$-, $`4312`$-, and $`3421`$-avoiding.
###### Proof.
The definition of the poset $`P(w)`$ and Theorem 6.4 indicate that $`P(w)`$ has a $`\widehat{1}`$ if and only if the union of any two decreasing subsequences that intersect at least twice is itself a decreasing subsequence.
Suppose there are decreasing subsequences in $`w`$ of lengths $`k_1,k_23`$ that intersect $`i2`$ times, for $`i<k_1,k_2`$. Let $`k=i+1`$, and choose a $`k+1`$ element subsequence of $`k_1\mathrm{}1k_2\mathrm{}1`$ that includes $`k_1\mathrm{}1k_2\mathrm{}1`$ and one more element from each descending subsequence. Let $`p๐_{k+1}`$ be the resulting pattern. No $`\widehat{1}`$ in $`P(w)`$ is equivalent to there being subsequences so that
$$p=(k+1)k\mathrm{}(j+2)j(j+1)(j1)\mathrm{}21$$
for some $`j[1,k]`$.
There are two ways to place a $`2k`$-gon in a zonotopal tiling of $`X(p)`$, but these overlapping $`2k`$-gons do not both lie in any larger centrally symmetric polygon. The permutation $`p`$ is always vexillary, so Theorem 3.8 implies that $`P(w)`$ will not have a $`\widehat{1}`$ if $`w`$ contains such a $`p`$.
Therefore, considering the permutation $`p`$ for each possible $`j`$, the poset $`P(w)`$ has a $`\widehat{1}`$ if and only if $`w`$ is $`4231`$-, $`4312`$-, and $`3421`$-avoiding. โ
The permutations for which $`P(w)`$ has a $`\widehat{1}`$ have recently been enumerated by Toufik Mansour in .
## 7. The Freely Braided Case
Although the graph $`G(w)`$ and poset $`P(w)`$ are not known in general, there is a class of permutations for which these objects can be completely described. This paper concludes with a study of this special case.
In and , Green and Losonczy introduce and study โfreely braidedโ elements in simply laced Coxeter groups. In the case of type $`A`$, these are as follows.
###### Definition 7.1.
A permutation $`w`$ is *freely braided* if every pair of distinct $`321`$-patterns in $`w`$ intersects at most once.
Equivalently, $`w`$ is freely braided if and only if $`w`$ is $`4321`$-, $`4231`$-, $`4312`$-, and $`3421`$-avoiding. The poset of a freely braided permutation has a unique maximal element by Theorem 6.13.
###### Example 7.2.
The permutation $`35214`$ is not freely braided because $`321`$ and $`521`$ are both occurrences of the pattern $`321`$, and they intersect twice. The permutation $`52143`$ is freely braided.
Mansour enumerates freely braided permutations in .
In , Green and Losonczy show that a freely braided $`w`$ with $`k`$ distinct $`321`$-patterns has
(6)
$$|C(w)|=2^k.$$
Moreover, in they show the following fact for any simply laced Coxeter group, here stated only for type $`A`$.
###### Proposition 7.3 (Green-Losonczy).
If a permutation $`w`$ is freely braided with $`k`$ distinct $`321`$-patterns, then there exists $`๐ขR(w)`$ with $`k`$ disjoint long braid moves.
###### Remark 7.4.
This means that there is a tiling of $`X(w)`$ with $`k`$ sub-hexagons, none of which overlap. Furthermore, equation (6) implies that flipping any sequence of these sub-hexagons does not yield any new sub-hexagons. Hence *every* tiling of $`X(w)`$ has exactly $`k`$ sub-hexagons, none of which overlap, and $`\widehat{1}`$ in $`P(w)`$ corresponds to the zonotopal tiling with rhombi and $`k`$ hexagons.
From Remark 7.4, the structures of the graph $`G(w)`$ and the poset $`P(w)`$ are clear for a freely braided permutation $`w๐_n`$.
###### Theorem 7.5.
Let $`w`$ be freely braided with $`k`$ distinct $`321`$-patterns. The graph $`G(w)`$ is the graph of the $`k`$-cube, and the poset $`P(w)`$ is isomorphic to the face lattice of the $`k`$-cube without its minimal element.
###### Example 7.6.
The permutation $`243196587`$ is freely braided. Its three $`321`$-patterns are $`431`$, $`965`$, and $`987`$. Figures 5 and 6 depict its graph and poset.
## 8. Acknowledgments
Particular thanks are due to Richard Stanley for his continued guidance and for the suggestion to study reduced decompositions. Anders Bjรถrner provided helpful advice and discussion, and together with Richard Stanley coordinated the semester on algebraic combinatorics at the Institut Mittag-Leffler, during which much of the research for this paper occurred. Thanks are also owed to John Stembridge for his referral to the work of Elnitsky and Green and Losonczy. Finally, the thoughtful suggestions of two anonymous referees have been greatly appreciated.
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# The VLT-UVES survey for molecular hydrogen in high-redshift damped Lyman-๐ผ systems: Physical conditions in the neutral gas
## 1 Introduction
Damped Ly-$`\alpha `$ (DLA) systems seen in QSO spectra are characterized by very large neutral hydrogen column densities: $`N(`$i$`)2\times 10^{20}`$ cm<sup>-2</sup>. Such an amount of neutral gas is usually measured through local spiral disks. The case for DLA systems to arise through proto-galactic disks is further supported by the fact that the cosmological density of the absorbing gas at $`z_{\mathrm{abs}}3`$ is of the same order of magnitude as the cosmological density of stars at present epochs (Wolfe 1995). Moreover, the presence of heavy elements ($`Z1/10Z_{}`$) suggests that DLAs are located in over-dense regions where star formation activity takes place (Pettini et al. 1997) and at low and intermediate redshifts strong metal line systems and DLAs have been demonstrated to be associated with galaxies (e.g. Bergeron & Boissรฉ, 1991; Le Brun et al. 1997). It has also been shown that the profiles of the lines arising in the neutral gas show evidence for rotation (e.g. Prochaska & Wolfe 1997). However, hydrodynamical simulations have shown that the high redshift progenitors of present-day galactic disks could look like an aggregate of well separated dense clumps. In fact, the kinematics seen in the absorption line profiles of DLAs could be explained by relative motions of the clumps with little rotation (Haehnelt et al. 1998; Ledoux et al. 1998).
Studying the star-formation activities in DLAs is very important for the understanding of galaxy formation in the Universe. Recently, Wolfe et al. (2003a, 2003b, 2004) have shown that, even if DLAs sustain only a moderate star-formation activity, they will contribute appreciably to the global star-formation rate (SFR) density at high redshifts. The SFR in DLAs can be estimated either by detecting the galaxies responsible for DLAs or by inferring the intensity of the UV field in DLAs using the induced excitation of atomic and molecular species. In the latter case, it is important to have a clear understanding of the physical conditions in the gas to derive an accurate estimate of the SFR. In the case of the Galactic ISM, rotational excitations of H<sub>2</sub> (see Browning et al. 2002 and references there in) and fine-structure excitations of C i, C ii, O i and Si ii are used to derive the physical state of the absorbing gas (see for example Welty et al. 1999). Detecting and studying these transitions in DLAs is the first step toward understanding the physical conditions and hence the star-formation activity in DLAs.
Molecular hydrogen is ubiquitous in the neutral phase of the interstellar medium (ISM) of galaxies. Formation of H<sub>2</sub> is expected on the surface of dust grains, if the gas is cool, dense and mostly neutral, and from the formation of H<sup>-</sup> ions if the gas is warm and dust-free (see e.g. Jenkins & Peimbert 1997; Cazaux & Tielens 2002). As the former process is most likely dominant in the neutral gas associated with DLA systems, it is possible to obtain an indirect indication of the dust content in DLAs without depending on extinction and/or heavy element depletion effects. Moreover, by determining the populations of different H<sub>2</sub> rotational levels, it is possible to constrain kinetic and rotational excitation temperatures and particle densities. Effective photo-dissociation of H<sub>2</sub> takes place in the energy range $`11.113.6`$ eV through Lyman- and Werner-band absorption lines and the intensity of the local UV radiation field can therefore be derived from the observed molecular fraction. A direct determination of the local radiation field could have important implications in bridging the link between DLA systems and star-formation activity at high redshifts.
We have searched for molecular hydrogen in DLA and sub-DLA systems at high redshift ($`z_{\mathrm{abs}}>1.8`$), using UVES at the VLT down to a detection limit of typically $`N(`$H$`{}_{2}{}^{})2\times 10^{14}`$ cm<sup>-2</sup> (see Ledoux et al. 2003). Out of the 33 systems in our sample, 8 have firm and 2 have tentative detections of associated H<sub>2</sub> absorption lines. In all of the systems, we measured metallicities relative to Solar, \[X/H\] (with either X$`=`$Zn, or S, or Si), and depletion factors of iron, \[X/Fe\], supposedly onto dust grains. Although H<sub>2</sub> molecules are detected in systems with depletion factor, \[Zn/Fe\], as low as 0.3, the systems where H<sub>2</sub> is detected are usually amongst those with the highest metallicities and depletion factors. In particular, H<sub>2</sub> is detected in the three systems with the largest depletion factors. Moreover, in two different systems, one of the H<sub>2</sub>-detected components has \[Zn/Fe$`]>1.5`$. This directly demonstrates that a large amount of dust is present in the components where H<sub>2</sub> is detected. The mean H<sub>2</sub> molecular fraction, $`f=2N(`$H$`{}_{2}{}^{})/[2N(`$H$`{}_{2}{}^{})+N(`$H i$`)]`$, in DLA systems is generally small (typically $`\mathrm{log}f<1`$) and similar to what is observed in the Magellanic Clouds. There is no correlation between the amount of molecules and the neutral hydrogen column density; in particular, two systems where H<sub>2</sub> is detected have $`\mathrm{log}N(`$H i$`)<20.3`$. Approximately 50 percent of the systems have $`\mathrm{log}f<6`$: this is probably a consequence of a reduced formation rate of H<sub>2</sub> onto dust grains (probably because the gas is warm, $`T>1000`$ K) and/or of an enhanced ionizing flux relative to what is observed in our Galaxy.
In this work, we present additional high S/N ratio data on three of the DLA systems in which H<sub>2</sub> is detected and the results of multi-component Voigt profile fits to neutral and singly ionized species (including C i, C i and C ii) in all the DLAs in our sample. We estimate the range of physical conditions in the neutral gas using standard techniques that are used in ISM studies. The paper is organized as follows. In Section 2, we give the details of the additional data and present the new fits to the H<sub>2</sub> absorption lines in the corresponding three systems. In Section 3, we discuss the relative populations of different H<sub>2</sub> rotational levels deriving information on the physical state of the gas by comparing the DLA observations with Galactic ISM, SMC, and LMC data. In Sections 4 and 5 we discuss, respectively, the fine-structure excitation of C i and the ionization state of Carbon. In Section 6 we study the C ii excitation in detail. Finally, we summarize our results and discuss various implications of the overall study in Section 7.
## 2 Data sample
The Ultraviolet and Visible Echelle Spectrograph (UVES; Dekker et al. 2000), installed at the ESO VLT 8.2-m telescope unit Kueyen on Mount Paranal in Chile was used to search for H<sub>2</sub> in a large sample of DLAs. The sample and data reduction procedure are described in detail in Ledoux et al. (2003). Observations and details of Voigt profile analysis of H<sub>2</sub> and metal line absorption lines toward Q 0013$``$004, Q 0551$``$366 and Q 1232$`+`$082, along the lines of sight of which H<sub>2</sub> is detected, are described in, respectively, Petitjean et al. (2002), Ledoux et al. (2002) and Srianand et al. (2000). The Voigt profile fits to H<sub>2</sub> and other metal lines at $`z_{\mathrm{abs}}`$= 3.024 toward Q $`0347383`$, $`z_{\mathrm{abs}}`$= 2.595 toward Q $`0405443`$ and $`z_{\mathrm{abs}}`$= 2.0868 toward Q $`1444+014`$ are discussed in Ledoux et al. (2003). For systems in which H<sub>2</sub> is not detected, Ledoux et al. (2003) have provided upper limits on $`N`$(H<sub>2</sub>) together with mean metallicities and depletion factors.
Recently, we have obtained additional higher spectral resolution spectra ($`R`$ 55,000) of Q 0347$``$383 and Q $`0405443`$ as a part of our ongoing programme on cosmic variation of the electron-to-proton mass ratio (Petitjean et al. 2004). Nine exposures of 1.5 h each were taken for each of the quasars over six nights under sub-arcsec seeing conditions in January 2002 and 2003 for Q 0347$``$383 and Q $`0405443`$, respectively. We have also obtained additional data of Q 1232$`+`$082 to study the HD lines that are detected in the DLA (Varshalovich et al. 2002). Spectra were reduced using the UVES pipeline and addition of individual exposures were performed using a sliding window and weighting the signal by the errors in each pixel. We detect a new H<sub>2</sub> component at $`z_{\mathrm{abs}}`$= 2.59486 toward Q 0405$``$443 in addition to the strong component reported in Ledoux et al. (2003). We also present Voigt profile fits to the H<sub>2</sub> lines in the $`z_{\mathrm{abs}}`$= 2.811 system toward Q $`0528250`$. The single H<sub>2</sub> component seen in the lower spectral resolution CASPEC spectrum (Srianand & Petitjean 1998) is resolved into two distinct components in our new UVES spectra. For both these systems the Voigt profile fits to the H<sub>2</sub> Lyman and Werner band absorption lines are shown in Fig. 1 and resulting parameters are summarised in Table 1. This Table also gives the results of Voigt profile fits to H<sub>2</sub> for the $`z_{\mathrm{abs}}`$= 2.33772 toward Q 1232$`+`$082.
The main purpose of this paper is to provide a detail account of C i, C ii and other metal lines in DLAs of our sample and extract physical conditions in conjunction with the H<sub>2</sub> content reported in Ledoux et al. (2003). For the $`z_{\mathrm{abs}}`$= 2.139 system toward Tol 1037$`+`$014 and the $`z_{\mathrm{abs}}`$= 3.350 system toward Q $`11171329`$, we use the results presented in Srianand & Petitjean (2001) and Pรฉroux et al. (2002) respectively. For the rest of the systems we give here the results of the multicomponent Voigt profile fits. For this we use a Voigt-profile fitting code that determines the best fitting parameters (column density, velocity dispersion and redshift) using $`\chi ^2`$ minimization techniques (Chand et al. 2004). We use the oscillator strengths compiled in Table 1 of Ledoux et al. (2003) for metal ions and those given by Morton & Dinerstein (1976) for H<sub>2</sub>. In this article, we measure metallicities relative to Solar, \[X/H$`]\mathrm{log}[N(`$X$`)/N(`$H$`)]\mathrm{log}[N(`$X$`)/N(`$H$`)]_{}`$, with either X$`=`$Zn, or S, or Si, and depletion factors of iron, \[X/Fe$`]\mathrm{log}[N(`$X$`)/N(`$Fe$`)]\mathrm{log}[N(`$X$`)/N(`$Fe$`)]_{}`$, adopting the Solar abundances from Savage & Sembach (1996).
## 3 Determination of physical parameters using H<sub>2</sub> level population
In this section, we estimate different physical parameters from the column densities of H<sub>2</sub> in different J rotational levels.
### 3.1 Kinetic temperature of the gas
It is a standard procedure, in ISM studies, to use the ortho-to-para ratio (OPR) to infer the kinetic temperature of the gas assuming local thermodynamic equilibrium, LTE (Tumlinson et al. 2002, and references there in). Indeed, recent numerical investigations suggest that the OPR is a good tracer of the kinetic temperature over large regions of the parameter space (Shaw et al. 2004). For completeness, we first review our understanding of the OPR and outline the method for deriving the kinetic temperature before applying the method to the data.
#### 3.1.1 General outline
As the interconversion between para and ortho states involves a spin flip, it is not allowed for processes involving an isolated molecule (i.e., radiative processes cannot induce interconversion). Ortho/para interconversion is only possible through (i) spin exchange induced by collisions with protons (with a rate coefficient in the range $`10^{10}10^9\mathrm{cm}^3\mathrm{s}^1`$; see Dalgarno, Black & Weisheit 1973; Flower & Watt 1984 and Gerlich 1990) or with hydrogen atoms (with a rate coefficient an order of magnitude less than that of protons; Mandy & Martin 1993; Tinรฉ et al. 1997) and (ii) reactions on the surface of dust grains (Le Bourlot 2000). In the case of local thermodynamic equilibrium (LTE),
$$\mathrm{OPR}_{\mathrm{LTE}}=3\frac{\underset{\mathrm{J}=\mathrm{odd}}{}(2\mathrm{J}+1)\mathrm{exp}[\mathrm{BJ}(\mathrm{J}+1)/\mathrm{T}]}{_{\mathrm{J}=\mathrm{even}}(2\mathrm{J}+1)\mathrm{exp}[\mathrm{BJ}(\mathrm{J}+1)/\mathrm{T}]}$$
(1)
where, J, is the rotational quantum number, $`B`$ is the rotational constant of H<sub>2</sub> ($`B`$ = 85.3 K), and $`T`$ is either the kinetic temperature of the gas (when OPR is governed by spin-exchange collisions) or the formation temperature (when OPR is governed by H<sub>2</sub> formation on the surface of dust grains with LTE distribution characterized by the formation temperature $`T_{\mathrm{form}}`$; see Sternberg & Neufeld 1999; Takahashi 2001). The equilibrium temperature, T(OPR), can be obtained using the observed value of OPR and Eq. 1. This will trace the kinetic temperature of the gas if spin-exchange collisions are mainly responsible for the observed OPR.
If the gas is dense and cold and if most of the H<sub>2</sub> molecules are in the J = 0 and J = 1 levels then,
$$\mathrm{OPR}_{\mathrm{LTE}}\frac{\mathrm{N}(\mathrm{J}=1)}{\mathrm{N}(\mathrm{J}=0)}=9\times \mathrm{exp}(170.5/\mathrm{T}_{01}).$$
(2)
In the case of very optically thick molecular gas for which there is enough self-shielding, the $`N`$(J=1)/$`N`$(J=0) ratio can be maintained at its Boltzmann value and the excitation temperature, $`T_{01}`$, equals the kinetic temperature. Savage et al. (1977) measured a mean excitation temperature, T<sub>01</sub> = 77$`\pm `$17 K, for the galactic ISM. This is consistent with the mean temperature of the ISM measured using 21 cm absorption lines. Thus it is widely believed that when there is sufficient shielding (i.e log $`N`$(H<sub>2</sub>) cm<sup>-2</sup> $``$ 16.5), $`T_{01}`$ is a reasonably good tracer of the kinetic temperature. This is because in the shielded region, H<sub>2</sub> photodissociation time-scale can be larger than the time-scale for charge exchange collision (Flower & Watt, 1984). Also a recent multi-wavelength study of Galactic sightlines show the T<sub>01</sub> measured in optically think cases closely follow the spin temperature measured from 21 cm observations (see Roy et al 2005).
The excitation temperature, $`T_{ij}`$, between different rotational levels (say J= i and j) of a given species (either ortho or para H<sub>2</sub>) can be obtained using,
$$\frac{N(J=j)}{N(J=i)}=\frac{2j+1}{2i+1}\mathrm{exp}(B[j(j+1)i(i+1)]/T_{ij}).$$
(3)
Unlike OPR, this ratio can be altered by radiation pumping and formation pumping in addition to collisions. If collisions dominate the rotational excitation then $`T_{ij}`$ will be equal to T(OPR). Presence of formation pumping and/or UV pumping will make $`T_{ij}>T(OPR)`$. In the following section we discuss various temperature estimates from the DLA sample.
#### 3.1.2 Kinetic temperature of the H<sub>2</sub> components
In our sample, H<sub>2</sub> is only detected in J$`5`$ levels of the vibrational ground state. Thus we compute the OPR by summing the H<sub>2</sub> column densities for levels with J $``$ 5. The observed value of the OPR for each DLA is given in column #8 of Table 2. We calculate T(OPR) from the measured OPR for individual systems using Eq. 1(see column #9 of Table 2). When the kinetic temperature (or formation temperature) is high (i.e., $`T200`$ K) the OPR reaches 3, the value expected based on spin statistics. For a kinetic temperature similar to that seen in the cold neutral medium of our Galaxy ($``$ 80 K) the expected OPR under LTE assumption is less than 1. From Table 2, it is clear that the LTE temperatures measured from the OPR for DLAs are most of the times higher than 80 K (the mean found in the Galactic ISM.)
In Fig. LABEL:fig2 we plot the observed values of the OPR against the total H<sub>2</sub> column density in the ISM (triangles), LMC (squares), SMC (asterisks) and DLAs (circles with error-bars). It is apparent that most of the OPR values in DLAs are significantly different from 3 (see also column #8 of Table 2). The distribution of the OPR as a function of $`N`$(H<sub>2</sub>) in DLAs is consistent with that observed along Galaxy, LMC and SMC sightlines (see Fig. LABEL:fig2) when $`\mathrm{log}`$ N(H<sub>2</sub>)$``$16 and for rest of the components OPR in DLAs are systematically higher than that measured in the Galaxy, LMC and SMC. For example, OPR $``$ 3 is seen only along sightlines with low H<sub>2</sub> optical depth (i.e $`N`$(H<sub>2</sub>$`10^{16}\mathrm{cm}^2`$) in the Galaxy, LMC and SMC (see Fig. LABEL:fig2). On the contrary, out of the three DLA components with OPR $`3`$, two, at $`z_{\mathrm{abs}}`$= 2.08680 and $`z_{\mathrm{abs}}`$= 2.08692 toward Q 1444+014, are optically thick in the Lyman band absorption lines.
We next investigate the dependence of $`T_{01}`$ (measured using Eq. 2) on the total H<sub>2</sub> column density (see Fig LABEL:figext). Individual values measured in DLAs are listed in Table 2 (see column #5). The large errors on both $`N`$(H<sub>2</sub>) and $`T_{01}`$ are mostly a consequence of the difficulty to measure the Doppler parameter when the lines are saturated. In the case of the $`z_{\mathrm{abs}}`$= 1.96685 component toward Q $`0013004`$ the uncertainty is a consequence of line blending (see Petitjean et al. 2002). The vertical dotted lines show the mean and 1$`\sigma `$ range of $`T_{01}`$ measured by Savage et al. (1977). The data points from the Magellanic clouds (Tumlinson et al. 2002) are consistent with this range (mean $`T_{01}=82\pm 21`$ K). As in the case of the OPR, most of the measurements from DLAs with optically thick H<sub>2</sub> (i.e log N(H<sub>2</sub>)$`16.5`$) are well separated from that of the ISM and Magellanic clouds (Fig. LABEL:figext) and the spread seen in the optically thin case is consistent with that seen in local ISM. Note that the system with lowest molecular content ($`z_{\mathrm{abs}}`$= 3.02489 toward $`0347383`$) has $`N`$(J=0) an order of magnitude lower than $`N`$(J =1). $`T_{01}`$ can not be computed in this case as the maximum expected column density ratio, $`N`$(J=1)/$`N`$(J=2), is 9 under LTE conditions (see Eq. 2). For the high optical depth clouds (i.e, log $`N`$(H<sub>2</sub>) cm<sup>-2</sup> $``$ 16.5) in DLAs the mean $`T_{01}`$ is 153$`\pm `$78 K. In most of the components the two temperatures $`T_{01}`$ and $`T`$(OPR) are consistent within errors. This is mainly because most of the H<sub>2</sub> molecules reside in the ground states.
In summary, if we assume LTE then $`T`$(OPR) and $`T_{01}`$ measured in DLAs (with log N(H<sub>2</sub>)$``$16.5) at high redshift are on an average higher than that measured in the ISM, LMC and SMC sightlines. In this high N(H<sub>2</sub>) range $`T_{01}`$ is expected to trace the kinetic temperature. However, in the case of optically thin systems the $`T`$(OPR) (or $`T_{01}`$) measured in DLAs are consistent with that measured in LMC, SMC and Galactic sightlines. Under the LTE assumption we find that H<sub>2</sub> components in DLAs have kinetic temperatures in the range 100$``$200 K.
### 3.2 Rotational excitation
The rotational level populations are affected by particle collisions, UV pumping, and formation pumping. While the collisional excitation plays a significant role in populating the low-J levels, those with J$``$3 are usually populated by formation processes and UV pumping. In what follows, we discuss the excitation of H<sub>2</sub> as seen in DLAs and compare with ISM, LMC and SMC sightlines.
#### 3.2.1 Low-J excitation
The collisional contribution to the excitation of H<sub>2</sub> can be investigated by studying the $`N`$(J=2)/$`N`$(J=0) and $`N`$(J=3)/$`N`$(J=1) ratios. In general J = 2 and J = 3 levels can also be populated by deexcitation of H<sub>2</sub> formed in the high-J states (usually referred to as formation pumping) or through UV pumping. The collisional excitation rate for the J=0$``$2 transition is about an order of magnitude higher than that of the J=1$``$3 transition for kinetic temperatures in the range 100 to 300 K (Forrey et al. 1997). The spontaneous decay rate from J = 3 is an order of magnitude smaller that from J = 2. This means that the ground and first excited states of para-H<sub>2</sub> can be thermalised at lower densities compared to that of ortho-H<sub>2</sub>. In Fig. LABEL:h2col, we plot as a function of temperature the critical hydrogen density for which the collisional deexcitation rate becomes equal to the spontaneous decay rate for the J=2$``$J=0 transition. It is clear from this figure that the hydrogen density has to be high (in the range 60$``$175 cm<sup>-3</sup>) in order for the $`N`$(J=2)/$`N`$(J=0) ratio to be equal to the LTE value corresponding to typical kinetic temperatures inferred from the OPR (i.e 100 to 300 K).
In Fig. LABEL:fig3 we plot the ratio, $`N`$(J=2)/$`N`$(J=0), observed in DLAs, the Galaxy, LMC and SMC as a function of the total H<sub>2</sub> column density. The vertical dotted lines in the figure shows the expected values of the ratio for four different excitation temperatures assuming LTE. Values of the excitation temperature $`T_{02}`$ for individual DLA H<sub>2</sub> components obtained using Eq. 3 are given in Table 2. The observed excitation temperatures are in the range 100 to 600 K with most of them at $`T_{02}150300`$ K. If the level populations are in LTE then the required hydrogen density to maintain the equilibrium is $`65150`$ cm<sup>-3</sup> (see Fig. LABEL:h2col). We can see from Fig. LABEL:fig3 that in DLAs where H<sub>2</sub> is optically thick, the $`N`$(J=2)/$`N`$(J=0) ratio is larger than that seen in similar gas of the Galactic ISM, LMC and SMC. It can be seen from Table 2 that, in most of the DLAs, $`T_{01}`$ is lower than or equal to $`T_{02}`$ (see Table 2). This is very much the case as well in most of the sightlines through the ISM and Magellanic clouds. It is well known that, due to a lower value of the Einstein coefficient of the J = 2 level compared to those of higher J levels, the UV and formation pumping processes can lead to enhancing the J = 2 level compared to the J = 0 level. Thus the higher values of $`N`$(J=2)/$`N`$(J=0) seen in DLAs can be explained by higher pressure in the gas and/or higher radiation field.
Fig. LABEL:fig4 gives the $`N`$(J=3)/$`N`$(J=1) ratio as a function of $`N`$(H<sub>2</sub>). The vertical dotted lines in the figure shows the expected value of the ratio for four different excitation temperatures under the LTE assumption. The measurements in DLAs are consistent with local measurements and the excitation temperature $`T_{03}`$ is in the range 100$``$680 K (see Table 2).
In Fig. LABEL:fig4a we plot the $`N`$(J=3)/$`N`$(J=1) ratio versus the $`N`$(J=2)/$`N`$(J=0) ratio. If formation and UV pumping contribute appreciably to populate the J = 2 and J = 3 levels then we expect a tight relationship between the two quantities. The dotted line in the figure gives the expected relationship between the ratios under LTE. In the case of sightlines through the Galactic ISM, the LMC or SMC, the $`N`$(J=3)/$`N`$(J=1) ratio is higher than what is expected from the $`N`$(J=2)/$`N`$(J=0) value under LTE (or, $`T_{13}`$ is higher than $`T_{02}`$). In the case of DLAs, most of the components have $`T_{13}`$ close to $`T_{02}`$ (points are on top of the dotted line). Note that these excitation temperatures are different from $`T_{01}`$. This clearly means that UV pumping and formation pumping are not negligible even for the excitation of the low J levels. The nature of the local radiation field can be probed using excitations of J$``$ 3 levels. This is what we do in the following Section.
#### 3.2.2 UV radiation field: High-J excitation
It is known that in the photodissociation regions (PDRs) the J = 4 and J = 5 rotational levels are populated predominantly by cascades following the formation of excited molecules and UV pumping from the low-J states. As radiative decay time-scales for these levels are very short compared to the collisional time-scales, spontaneous decay is the main deexitation process. Among the two populating processes the UV pumping is an optical depth dependent process while the formation pumping is independent of optical depth. In an optically thick cloud, UV pumping is efficient in a thin shell surrounding the cloud. In the interior of the cloud, UV pumping becomes important only when the column density becomes very large (i.e. through absorption in the damping wings).
In Fig. LABEL:fig6 we plot log $`N`$(H<sub>2</sub>) as a function of log $`N`$(J=4)/$`N`$(J=0) as measured in DLAs and along the ISM, LMC and SMC sightlines. As expected, a strong anti-correlation is present in the data, including DLAs. H<sub>2</sub> absorption lines in DLA systems have no strong overlapping wings. Therefore the high-J excitation is mostly due to photo-absorption in the systems with log $`N`$(H<sub>2</sub>) $``$16.5 and to H<sub>2</sub> formation in the systems with higher column densities. Following analytic prescription by Jura (1975) we can write,
$$p_{4,0}\beta \left(0\right)n\left(H_2,J=0\right)+0.24Rn\left(H\right)n=A\left(42\right)n\left(H_2,J=4\right)$$
(4)
Here, $`\beta (0)`$, $`p_{4,0}`$ are, respectively, the photo-absorption rate in the Lyman and Werner bands and the pumping efficiency from J = 0 to J = 4; $`\mathrm{A}(42)`$ is the spontaneous transition probability between J = 4 and J = 2 and $`R`$ is the formation rate of H<sub>2</sub>. Neglecting the second term in the left hand side of Eq. 4 leads to a conservative upper limit on the UV radiation field. The vertical dashed lines in Fig. LABEL:fig6 represent the corresponding predicted values of the $`N`$(J=4)/$`N`$(J=1) ratio for $`\beta (0)=2\times 10^{10}`$s<sup>-1</sup> (that is, approximatively the mean radiation field in the ISM) and $`\beta (0)=2\times 10^9`$s<sup>-1</sup>. It can be seen from Fig. LABEL:fig6 that for log $`N`$(H<sub>2</sub>) less than 16.5 the $`N`$(J=4)/$`N`$(J=0) ratio in DLAs is of the order of or slightly higher than that seen in the ISM of our Galaxy. Quantitatively the upper limits in most of the systems are consistent with $`2\times 10^{10}\beta (0)2\times 10^9`$s<sup>-1</sup>. This probably means the optically thin H<sub>2</sub> components without detectable H<sub>2</sub> absorption lines from the J = 4 state arise in gas embedded in a UV field with intensity similar to (or slightly higher than) that of the mean ISM field.
There are two optically thin components in our sample ($`z_{\mathrm{abs}}`$= 1.96822 toward Q 0013$``$004 and 3.02489 toward Q 0347$``$383) that show detectable J = 4 H<sub>2</sub> absorption lines. Detailed analysis of these component suggests an ambient field intensity consistent with few times the mean ISM field intensity (Petitjean et al. 2002; Levshakov et al. 2002). The same conclusion was derived by Reimers et al. (2003) for the optically thin H<sub>2</sub> component at $`z_{\mathrm{abs}}`$= 1.15 system toward HE $`05154414`$.
The above ratio has similar values at high log $`N`$(H<sub>2</sub>) in DLAs and in our Galaxy. This is a hint for the formation pumping in DLAs with high $`N`$(H<sub>2</sub>) being similar to the local one. There are two optically thick components (at $`z_{\mathrm{abs}}`$= 2.59471 toward Q 0405$``$443 and $`z_{\mathrm{abs}}`$= 2.08696 toward Q 1444$`+`$014) that do not show detectable absorption lines from the J = 4 state. In the latter system the ratio $`N`$(J=4)/$`N`$(J=0)$`10^4`$. This is much lower than the values seen in the ISM at similar total $`N`$(H<sub>2</sub>) and could be a consequence of lower H<sub>2</sub> formation rate in this system. High values of the radiation field intensity were inferred for some of the optically thick components when the contribution of the second term in Eq. 4 is estimated using the average metallicity and dust depletion (Ge & Bechtold 1997; Petitjean et al. 2000; Ge, Bechtold & Kulkarni 2001).
## 4 Analysis of Carbon absorption lines
As the ionization potential of C i is 11.2 eV, the ionization state of Carbon is sensitive to the same photons that destroy H<sub>2</sub>. Therefore, C i is usually a good tracer of the physical conditions in the molecular gas (see however Srianand & Petitjean 1998). In what follows we investigate the relationship between the detectability of C i absorption line and other measurable quantities in our spectra. We derive additional constraints on the physical conditions in DLAs using C i fine-structure absorption lines.
### 4.1 Detectability of C i absorption lines
The results of simultaneous Voigt profile fitting to C i, C i and C i<sup>โโ</sup> absorption lines in our sample are summarised in Table 3.
In the interstellar medium of our Galaxy, all clouds with log $`N`$(H i$`21`$ have log $`N`$(H<sub>2</sub>$`>19`$ and log $`N`$(C i$`>14`$ (Jenkins & Shaya 1979; Jenkins, Jura & Loewenstein 1983). In our sample, C i absorption lines are detected in most of the DLAs that show H<sub>2</sub> absorption lines(see also Ge & Bechtold 1999). There are three exceptions: the components at $`z_{\mathrm{abs}}`$= 2.59471 and 2.59486 toward Q 0405$``$443 and at $`z_{\mathrm{abs}}`$= 2.81100 toward PKS 0528$``$250. These C i non-detections are surprising as the H<sub>2</sub> absorption lines from these components are optically thick so that C i is expected to be conspicuous.
Usually, DLAs in which no H<sub>2</sub> is detected through the whole profile do not show any detectable C i absorption line (with a typical upper limit of 10<sup>12</sup> cm<sup>-2</sup>). The only exception is the high-metallicity sub-DLA at $`z_{\mathrm{abs}}`$= 2.139 toward Tol 1037$``$270 (see Srianand & Petitjean 2001). On the contrary, in DLAs where H<sub>2</sub> is detected, some components show detectable C i absorption line without detectable molecular absorption ($`N`$(H<sub>2</sub>)$`10^{14}`$cm<sup>-2</sup>). This is the case in Q 0013$``$004 (Petitjean et al. 2002) and Q 0551$``$366 (Ledoux et al. 2002).
Note that C i is also detected at $`z_{\mathrm{abs}}`$= 2.28749 toward Q 2332$``$094 but the presence of H<sub>2</sub> molecules can not be probed in this system due to the presence of an intervening Lyman limit system. The sub-DLA at $`z`$ = 1.15 toward HE 0515$``$4414 shows C i and H<sub>2</sub> absorption lines (Quast et al. 2002; Reimers et al. 2003). C i absorption lines have also been detected at $`z_{\mathrm{abs}}`$= 1.776 system toward Q 1331+170 (Chaffee et al. 1988). Presence of H<sub>2</sub> in this system is recently reported (Cui et al., 2004).
#### 4.1.1 Dependence on H<sub>2</sub> column density
Jenkins & Shaya (1979) found $`N`$(C i) does not scale linearly with either of $`N`$(H i), $`N`$(H<sub>2</sub>) or $`N`$(H<sub>total</sub>) in the Galactic ISM. They explained this behavior as a result of strong differences in the response of C i, H i and H<sub>2</sub> to physical conditions (electron density, temperature etcโฆ), coupled with marked variations of these conditions from one cloud to the other. In Fig. LABEL:c1fig1, we plot the C i column density as a function of H<sub>2</sub> column density in individual components. Among the systems that show H<sub>2</sub> absorption lines(filled circles with error-bars) there is no clear trend between $`N`$(H<sub>2</sub>) and $`N`$(C i) even though the presence of C i absorption lines usually indicate the presence of H<sub>2</sub> (see discussion above).
#### 4.1.2 The Carbon ionization state
The probability of detecting C i is expected to be higher in systems with higher $`N`$(H i) and/or metallicity. Ideally, we would like therefore to know $`N`$(H i) for each individual C i components. This is not possible as all components are definitely blended in one strong H i DLA absorption line. Estimation of N(H i) is possible when the H<sub>2</sub> component is well separated from the rest of the components (as in $`z_{\mathrm{abs}}`$= 1.96822 toward Q 0013$``$004) or when 21 cm observations are available (as in the case of H<sub>2</sub> components toward Q 0528$``$250). Note that the presence of very strong C i absorption line in the component at $`z_{\mathrm{abs}}`$= 1.96822 toward Q 0013$``$004 (that has log N(H i)$``$19.4) is mainly due to high metallicity (Petitjean et al. 2002). Whereas the absence of C i in the component at $`z_{\mathrm{abs}}`$= 2.81100 and the weakness of C i line of the component at $`z_{\mathrm{abs}}`$= 2.81112 component toward Q 0528$``$250 are probably due to excess radiation field from the QSO.
The ionization state of Carbon is difficult to determine as the C ii$`\lambda 1334`$ absorption line is usually highly saturated. We can however partly overcome this difficulty assuming that conditions are fairly homogeneous in the DLA system. Under the assumptions that the enrichment of Carbon follows that of $`\alpha `$elements and that the relative depletion between Sulfur and Carbon is negligible, we can use the well determined $`N`$(S ii) as an indicator of $`N`$ (C ii) (see Fig. LABEL:c1fig2). In the sun, the Carbon abundance is 1.28 dex higher than that of Sulfur and typical depletion of Carbon relative to Sulfur in the Cold ISM is 0.4 dex (see Table 5 of Welty et al. 1999). The two dotted lines in Fig. LABEL:c1fig2 give the expected correlation for log $`N`$(C i)/$`N`$(C ii) = $`3`$ (lower line) and $`2`$ (upper line) respectively, when relative solar abundances are used and it is assumed that there is no depletion of Carbon relative to Sulfur. The short-dashed lines give the same correlations when a depletion of Carbon relative to Sulfur of 0.4 dex is further assumed. If the absorbing gas originates from the CNM then log $`N`$(C i)/$`N`$(C ii) is expected to be more than $`3`$ (see Fig. 3 of Liszt 2002). Therefore, within uncertainties due to depletion, it is apparent from Fig. LABEL:c1fig2 that the DLA components with C i detections have a ionization state consistent with them originating from the CNM.
It is interesting to note that the distribution of $`N`$(S ii) is somewhat similar for components with both H<sub>2</sub> and C i absorption lines (filled circles), and for components with C i but no H<sub>2</sub> absorption lines (open circles). However C i column densities are typically lower in the components without H<sub>2</sub> suggesting, as expected, that when H<sub>2</sub> is seen, the C i/C ii ratio is larger.
Most of the upper limits on C i are consistent with $`N`$(C i)/$`N`$(C ii$``$ $`3`$ (see Fig. LABEL:c1fig2). This can mean either that the relative depletion of Carbon compared to Sulfur is larger than 0.4 dex in the CNM, which is unlikely, or that most of the DLA systems originate from the warm neutral medium (WNM) or warm ionized medium (WIM) where the above ratio can be as low as $`10^4`$.
#### 4.1.3 The effect of dust
In Fig. 11 we plot log $`N`$(C i) against the depletion factor defined as log ($`N`$(Fe ii)/$`N`$(X ii))$``$\[Fe/X\] with either X = Zn, S or Si. C i absorption line is not detected in systems with low depletion factors (i.e., \[Fe/Zn\] less than $``$0.5 dex) whereas components with higher depletion factors readily show detectable C i absorption lines. This trend is not surprising as there is a 4$`\sigma `$ correlation between the depletion factor and the metallicity of the gas in our sample (see Figure 12 of Ledoux et al. 2003). Depletion factors lower than 0.5 dex are usually seen in systems with \[Zn/H\] $`1`$. High depletion factor in high metallicity gas implies high dust content and hence high dust optical depth to the UV radiation. The absence of C i in components with low dust depletion is a combination of low metallicity and low dust optical depth to the UV radiation. It is worth remembering that similar relation exists between the detectability of H<sub>2</sub> and depletion (Fig. 14 of Ledoux et al. 2003).
### 4.2 i fine-structure excitation
In most of the DLAs with C i detections we also detect absorption lines from the excited fine-structure levels. It is therefore possible to use the relative populations of the C i ground state levels to discuss the particle density, the ambient UV radiation field and the temperature of the cosmic-microwave background radiation (see Bahcall et al., 1973; Meyer et al., 1986; Songaila et al. 1994; Ge, Bechtold & Black, 1997; Roth & Bauer, 1999; Srianand et al. 2000; Silva & Viegas, 2002; Quast et al. 2002). In the Galactic ISM, fine-structure excitation of C i has been used to study the distribution of thermal pressure (see Jenkins & Tripp 2001).
In Fig. LABEL:figc1 we plot the ratio $`N`$(C i)/$`N`$(C i) as a function of $`N`$(C i). It is clear that the C i column densities in DLAs (filled circles with error-bars) are at least an order of magnitude less than that measured in the ISM (stars). This is probably a consequence of lower metallicities and/or low H i content in DLA components. The important point is that the $`N`$(C i)/$`N`$(C i) ratio measured in DLAs is remarkably larger than in the Galaxy. However, while comparing the ISM and DLAs, it is important to remember that most of the sightlines used by Jenkins & Tripp (2001) have H<sub>2</sub> fraction orders of magnitude higher than what we measure in DLA components. As H<sub>2</sub> collisions are less efficient in populating the excited fine-structure state of C i we expect that for a given total hydrogen density (and a given kinetic temperature) $`N`$(C i)/$`N`$(C i) be higher in DLAs.
The horizontal dotted line in Fig. LABEL:figc1 indicates the expected value of the ratio if it is assumed that the excitation is due to the CMBR only with a temperature $`T_{\mathrm{CMBR}}`$ = 8.1 K as expected at $`z`$ = 2, the typical redshift of our sample. It is clear that the CMBR field expected from the Big-Bang is not sufficient to explain the observed ratios and an extra contribution is required from collisional processes and/or the UV flux.
Collisions with H, He, e, p and H<sub>2</sub> can populate the excited state of C i. As H i is the dominant form of hydrogen in the gas, the contribution to the fine-structure excitation by H<sub>2</sub> collisions can be neglected. The electron and proton densities are expected to be very small, at least smaller by two orders of magnitude than the hydrogen density, and their contribution is also negligible (Keenan et al. 1986). The He i collisional rates are much less than that of H i and the He i/H i ratio is small which makes collisions with He i unimportant (see Fig. 1 of Silva & Viegas 2002). Thus in our analysis of the C i excitation we consider only collisions by neutral hydrogen. The rates are taken from Launay & Roueff (1977). The spontaneous decay rates are $`\mathrm{A}_{10}=7.93\times 10^8\mathrm{s}^1`$ and $`\mathrm{A}_{21}=2.68\times 10^7\mathrm{s}^1`$ (Bahcall & Wolf 1968). The corresponding CMBR excitation rate is derived from these values. The UV pumping rate in the cloud depends on the nature and strength of the UV radiation field. We assume that the UV intensity is the same as in the ISM of our Galaxy as suggested by the high$`J`$ excitation of H<sub>2</sub>.
The dashed lines in Fig. LABEL:figc1 give the expected ratio for, $`T_{\mathrm{CMBR}}`$ = 8 K, a UV radiation field like in the Galaxy (with an excitation rate of $`\mathrm{\Gamma }_{01}=7.55\times 10^{10}\mathrm{s}^1`$), hydrogen density in the range $`n_\mathrm{H}`$ = 20$``$250 cm<sup>-3</sup>, versus the kinetic temperature, $`T_{\mathrm{kin}}`$, in the range $`80200`$ K (consistent with $`T_{01}`$ measured in DLAs). It can be seen that a typical density range consistent with most of the observed point is, 20$`n_\mathrm{H}(\mathrm{cm}^3)`$150. The density ($`n_\mathrm{H}`$) and pressure ($`p/K`$) derived for individual components are summarized in column 9 and 10 respectively in Table. 3. Here, we assume $`T_{\mathrm{CMBR}}`$ = 2.7$`\times `$(1+$`z_{\mathrm{abs}}`$) and $`T`$ = $`T_{01}`$ in the case of H<sub>2</sub> detection and $`T`$ = 100 K otherwise. The derived pressure range in DLA components are higher than that typically measured in the galactic ISM (see Jenkins & Tripp 2001) and consistent with what is expected in a the cold neutral medium (CNM) with lower metallicity (Z$``$0.1 Z) and dust depletion (see Wolfire et al. (1995, 2001), Wolfe et al., (2003a, b) Srianand et al. 2005). Interestingly the derived density range in most of the components is close to the critical density for thermalising the H<sub>2</sub> $`N`$(J=2)/$`N`$(J=0) ratio (see Fig. LABEL:h2col).
## 5 Carbon ionization state
In the ISM, when hydrogen is neutral, most of the carbon is in the form of C ii so that we can use the C i/C ii ratio to derive the physical conditions in the gas. As already mentioned, the C ii absorption features are always badly saturated and we rely on the realistic assumption that C ii in the neutral phase can be traced by S ii and/or Si ii so that we can use the weak lines of these species to derive the C ii column densities in the components of interest.
Assuming photoionization equilibrium between C i and C ii and using the atomic data from Shull & Van Steenberg (1982) and Pรฉquignot et al. (1991), we can write,
$$\frac{n_\mathrm{e}}{\mathrm{\Gamma }}=4.35\times 10^{11}\frac{N(\mathrm{C}\mathrm{I})}{N(\mathrm{C}\mathrm{II})}\left(\frac{T}{10^4}\right)^{0.64}$$
(5)
Here, $`\mathrm{\Gamma }`$ is the photoionization rate for C i. In the local ISM, $`\mathrm{\Gamma }_{\mathrm{gal}}23.3\times 10^{10}`$ s<sup>-1</sup> (Pรฉquignot & Aldrovandi 1986). Here, we neglect the ion-molecular interaction and charge exchange reactions that may produce C i. Thus $`n_\mathrm{e}/\mathrm{\Gamma }`$ can be constrained once the temperature of the gas is known. Application of Eq. 5 to the cold Galactic ISM (with $`T`$ = 100 K) has resulted in $`n_\mathrm{e}0.14\pm 0.07`$ cm<sup>-3</sup> (Welty et al. 2002 and references there in). In the stable CNM considered by Wolfire et al. (1995) the electron density is expected to be in the range 0.01$`n_\mathrm{e}(\mathrm{cm}^3)`$0.02 for $`Z`$ = 0.1 $`Z_{}`$ and dust abundance one tenth of the Galactic ISM.
### 5.1 Systems with C i detections
First we concentrate on the systems with C i detections. We estimate the electron density assuming a UV field similar to the Galactic mean field (i.e., $`\mathrm{\Gamma }_{\mathrm{gal}}=2.5\times 10^{10}`$ s<sup>-1</sup>), $`T=T_{01}`$ for the H<sub>2</sub> components and $`T`$ =100 K otherwise. Individual values of $`n_\mathrm{e}`$ derived for these systems are given in Table 4. The electron density is in the range 0.7$`\times 10^2n_\mathrm{e}(\mathrm{cm}^3)4.9\times 10^2`$. Together with n<sub>H</sub> given in Table 3, this suggests that $`n_\mathrm{e}/n_\mathrm{H}10^3`$ for most of the systems. Therefore, the ionization state of the gas with H<sub>2</sub> and C i is similar to that in the CNM in a moderate radiation field.
### 5.2 Systems without H<sub>2</sub> and C i
In the case of systems in which neither H<sub>2</sub> nor C i are detected, we assume $`T`$ = 100 K and $`\mathrm{\Gamma }_{\mathrm{gal}}=2.5\times 10^{10}`$ s<sup>-1</sup> and obtain upper limits on the electron density. The results are plotted in Fig. LABEL:eden. We notice that the inferred electron densities are much smaller than in systems in which C i (and H<sub>2</sub>) are detected. The difference is even larger if we use only the systems with log $`N`$(C ii)$``$15. For this gas the inferred electron densities are less than 10<sup>-2</sup> cm<sup>-3</sup> (with a median of 10<sup>-3</sup> cm<sup>-3</sup>). This may indicate that the absorption originates from warm neutral medium (say $`T`$ = 8000 K). If the average radiation field in DLAs is similar to that in the Galactic ISM then the absence of C i in most of the DLAs could just be a consequence of lower densities (and/or higher $`T`$) in these systems. One can derive an independent estimate of the particle density using the excitation of C ii fine structure levels. This is what we do in the following Section.
## 6 Excited fine-structure line of C ii
### 6.1 Method to derive physical parameters
Under LTE, the column density ratio $`N`$(C ii)/$`N`$(C ii), can be written as,
$$\frac{N(\mathrm{C}\mathrm{II}^{})}{N(\mathrm{C}\mathrm{II})}=\frac{Q_{12}(\mathrm{e})\mathrm{n}_\mathrm{e}+\mathrm{Q}_{12}(\mathrm{H})\mathrm{n}_\mathrm{H}+\mathrm{\Gamma }_{12}(\mathrm{CMB})}{\mathrm{A}_{21}}$$
(6)
where, $`Q_{12}(\mathrm{e})=7.8\times 10^6exp[91.27/T]T^{0.5}`$ cm<sup>-3</sup> s<sup>-1</sup> and $`Q_{12}(\mathrm{H})=1.3\times 10^9exp[91.27/T]`$ cm<sup>-3</sup> s<sup>-1</sup> are the collisional excitation rates per unit volume for electrons and hydrogen atoms (Bahcall & Wolf 1968) with $`T`$ being the kinetic temperature of the gas. The Einsteinโs coefficient is $`A_{21}=2.291\times 10^6`$ s<sup>-1</sup>. The CMB pumping rate, $`\mathrm{\Gamma }_{12}(\mathrm{CMB})`$, equals 6.6$`\times 10^{11}`$ s<sup>-1</sup> and 1.1$`\times 10^9`$ s<sup>-1</sup> for, respectively, redshifts 2 and 3 (Silva & Viegas 2002). For a given temperature, the collisional excitation rate for electrons is orders of magnitude larger than that for hydrogen atoms. For example, when $`T`$ = 1000 K, whenever $`n_\mathrm{e}/n_\mathrm{H}`$ is larger than 5$`\times 10^3`$, collisions with electrons is the dominant process. UV pumping is an additional possible excitation mechanism. For the mean radiation field in our Galaxy, the UV pumping rate is 9.3$`\times 10^{11}`$ s<sup>-1</sup> (Silva & Viegas 2002). This is similar to or slightly lower than the CMB pumping rate for the range of redshift we consider in this study.
We compute the expected value of log $`N`$(C ii)/$`N`$(C ii) ratio under different situations and in particular the WNM and CNM solutions given in Table 3 of Wolfire et al. (1995). Note that if CMBR pumping alone is responsible for the excitation, the expected ratios are $`4.54`$ and $`3.31`$ for, respectively, $`z_{\mathrm{abs}}`$= 2 and 3. For the standard ISM (with stable pressure in the range 990$``$3600 cm<sup>-3</sup> K), we derive $`2.62`$log $`N`$(C ii)/$`N`$(C ii)$`2.20`$ for the CNM and $`3.65`$log $`N`$(C ii)/$`N`$(C ii)$`3.17`$ for the WNM. As the DLA gas has low metallicity and low dust content, the expected pressure should be higher in the two DLA phases (see Liszt 2002 and Wolfe et al. 2003). If we assume $`Z=0.1Z_{}`$ and dust to gas ratio one tenth of the ISM value (which is typical of DLAs) then we expect $`2.26`$log $`N`$(C ii)/$`N`$(C ii)$`2.08`$ for the CNM and $`3.39`$log $`N`$(C ii)/$`N`$(C ii)$`2.70`$ for the WNM in DLAs. Thus if DLAs originate from H i gas in a two-phase equilibrium, we expect the CMBR pumping to be sub-dominant compared to collisional excitation. In fact if the gas is completely neutral then from Eqs. (5) and (6) we derive,
$`n_\mathrm{H}`$ $`=`$ $`{\displaystyle \frac{\mathrm{A}_{21}\mathrm{N}(\mathrm{C}\mathrm{II}^{})}{\sigma _\mathrm{H}\mathrm{N}(\mathrm{C}\mathrm{II})}}{\displaystyle \frac{\sigma _\mathrm{e}\mathrm{n}_\mathrm{e}}{\sigma _\mathrm{H}}}`$ (7)
$`=`$ $`{\displaystyle \frac{\mathrm{A}_{21}\mathrm{N}(\mathrm{C}\mathrm{II}^{})}{\sigma _\mathrm{H}\mathrm{N}(\mathrm{C}\mathrm{II})}}{\displaystyle \frac{\sigma _\mathrm{e}\mathrm{N}(\mathrm{C}\mathrm{I})\mathrm{\Gamma }}{\sigma _\mathrm{H}\alpha _\mathrm{r}\mathrm{N}(\mathrm{C}\mathrm{II})}}.`$
This gives an independent estimate of n<sub>H</sub> which, by comparison with the estimate derived from the C i excitation, can lead to constraints on the radiation field. However, the ionization fraction (i.e $`n_\mathrm{e}/n_\mathrm{H}`$) and the temperature of the neutral absorbing gas must be accurately determined before densities can be derived using the $`N`$(C ii)/$`N`$(C ii) ratio. Indeed, the fact that, in DLAs, the Al iii absorption profile is very similar to that of neutral or singly ionized species has been used as evidence for the presence of ionized gas being mixed with the neutral gas in DLAs (Lu et al. 1996; Prochaska & Wolfe 1999; Howk & Sembach 1999; Wolfe & Prochaska 2000; Vladilo et al. 2001; Izotov et al. 2001). In the ionized gas (i) $`n_\mathrm{e}`$ as well as $`T`$ will be higher than what is expected in the warm or cold neutral gas and (ii) $`N`$(Si ii) will under-predict $`N`$(C ii) as the ionization corrections are different for the two species (see Fig. 1 of Izotov et al. 2001). Neglecting the presence of ionized gas can artificially enhance the derived n<sub>H</sub> values.
### 6.2 Frequency of C ii detection
In our sample, all the systems that show H<sub>2</sub> also show detectable C ii absorption line. The details of the fits to the C ii absorption line for these systems are summarized in Table 5 and shown in Fig. LABEL:vpc2star. The components toward Q $`0013004`$ and Q $`0528250`$ are badly blended and it is therefore not possible to fit $`N`$(C ii). We detect C ii in seven out of the 21 DLAs that do not show H<sub>2</sub> absorption lines. If we also include the 8 DLAs that show H<sub>2</sub>, about 50% (15 out of 29) of the DLAs in our sample show detectable C ii absorption line which is consistent with the finding by Wolfe et al. (2003).
In Fig. LABEL:c2sdet we plot the average (over the whole profile) $`N`$(C ii)/$`N`$(C ii) ratio measured in DLAs in our sample together with the measurements by Wolfe et al. (2003a, 2003b) against the total H i column density and the silicon metallicity. Here we use the total column density summed over all the components. $`N`$(C ii) is computed from $`N`$(Si ii) assuming solar abundance ratio without any ionization correction. From the upper panel in Fig. LABEL:c2sdet it can be seen that C ii is detected in all the systems with log $`N`$(H i) $``$21.0. Interestingly no such relationship exists between C i (or H<sub>2</sub>) and H i. Most the systems with log $`N`$(H i) $``$21.0 have $`N`$(C ii)/$`N`$(C ii) consistent with what is expected in CNM. On the contrary, the measured values of $`N`$(C ii)/$`N`$(C ii) in systems with lower $`N`$(H i) spread over more than an order of magnitude covering the expected ranges for WNM and CNM.
From the bottom panel, it can be seen that C ii is frequently detected in gas with high metallicity as already noticed by Wolfe et al. (2003a). Most of the systems that show C ii absorption line with lower $`N`$(H i) do have statistically higher metallicity. In the whole sample the number of systems with C ii detections that are consistent with CNM and WNM are approximately equal. Most of the upper limits on the ratio, measured in the metallicity range $`2.0Z_{}1.5`$, are lower than what would be expected from CNM gas and are consistent with WNM (or low density) gas. Interestingly these upper limits are lower than that seen in high latitude Galactic sightlines that are believed to be predominantly WNM gas. This means that the electron density (and therefore probably the total particle density) in these DLAs is probably quite small.
### 6.3 Systems with H<sub>2</sub> detection
Systems with H<sub>2</sub> detections (marked as filled circles in Fig. LABEL:c2sdet) have $`N`$(C ii)/$`N`$(C ii) consistent with CNM. We compute the allowed range of n<sub>H</sub> in these components using $`N`$(C ii), $`N`$(S ii), $`n_\mathrm{e}`$ from the C i excitation and $`T=T_{01}`$ if available or $`T`$ = 100 K (see Eq.7). The results are summarised in Column 6 of Table 4. This Table also gives upper limits on n<sub>H</sub> for components without H<sub>2</sub> in systems that show H<sub>2</sub>. It is to be remembered that we assume \[C/S\] in DLAs is \[C/S\]. Realistically Carbon can be under-abundant by up to a factor of 2. In that case the density will be higher than what we quote in the table. From Table 3 and 4 it is clear that for the H<sub>2</sub> components toward Q 0347$``$383, Q 0551$``$366 and Q 1232$`+`$082 the value of n<sub>H</sub> derived from both methods agree well. Such a comparison is not possible for the components toward Q 0013$``$004 and Q 0405$``$383 as C ii is blended in the former case and C i is not detected in the latter case. In the case of Q 1444+014 the derived hydrogen density based on C ii is lower than that derived using C i fine-structure excitation. However, in this system, Ledoux et al. (2003) have found a 5 km/s shift between the C i absorption line and that of singly ionized species. In addition these components showing relative depletion of Si with respect to S, it is possible that we have over estimated $`N`$(C ii). In summary, the $`n_\mathrm{H}`$ estimates based on the two methods are approximately consistent with one another. The excitation of the fine-structure levels of C i and C ii in the components with H<sub>2</sub> detection are consistent with high density and low temperature CNM gas.
### 6.4 Systems without H<sub>2</sub> detection
In this Section we focus our attention on the 7 systems in our sample that show C ii without H<sub>2</sub>. These systems do not show detectable C i absorption lines except the high-metallicity system at $`z_{\mathrm{abs}}`$ = 2.1391 toward Tol 1037$``$270. Apart from the $`z_{\mathrm{abs}}`$= 1.943 system toward Q 1157+014 that show 21 cm absorption line (Wolfe et al. 1981) there is no independent constraint on $`T`$ and $`n_\mathrm{H}`$. The identification of C ii at $`z_{\mathrm{abs}}`$= 2.422 toward Q $`0112+029`$ and $`z_{\mathrm{abs}}`$= 2.799 toward Q $`0135273`$ is based on C ii$`\lambda 1037`$ absorption line. These lines are well inside the Lyman-$`\alpha `$ forest and possible contamination by intervening H i absorption cannot be ruled out. For the rest of the systems, the identification and estimation of the C ii column density are secure.
We detect Al iii absorption lines in the 6 (out of 7) systems for which our spectra cover the expected wavelength range of the redshifted Al iii transitions. We estimate the fraction of Al in Al iii using the observed metallicity and the observed $`N`$(Al iii), assuming no depletion and solar relative abundances. Results are summarised in Table 6. From the Table, it can be seen that 1$``$7 % of Al is twice ionized. Using photoionization models from โCLOUDYโ we derive a typical ionization parameter $`3`$ log $`U`$ $`2`$ if the gas originates from a single slab irradiated by stellar spectrum with an effective black-body temperature of 30,000$``$40,000 K (also see Fig. 1 in Izotov et al. 2001). This implies that the average $`n_\mathrm{e}/n_\mathrm{H}`$ ratio along the line of sight is typically in the range 0.3 to 0.9. Thus there are enough electrons in the cloud so that collisions with electrons are dominant in the excitation of C ii.
The average density, $`n_\mathrm{H}`$, is derived using Eq. 7 and assuming three possible combinations of $`T`$ and $`n_\mathrm{e}/n_\mathrm{H}`$. The results are summarised in Table 6. When we use no additional constraints, the C ii observations alone are consistent with the gas having high density and low temperature (see Column 4 of Table 6).
The last column in the table gives the upper limit on $`n_\mathrm{H}`$ that will keep the equilibrium abundance of H<sub>2</sub> below our detection limit (i.e $`N`$(H<sub>2</sub>)$`10^{14}`$ cm<sup>-2</sup>). This value is computed using simple formation equilibrium of optically thin H<sub>2</sub> (Jura 1975)
$$n_\mathrm{H}=\frac{0.11\beta (0)N(\mathrm{H}_2)}{RN(\mathrm{HI})}$$
(8)
with, $`R`$ and $`\beta (0)`$, respectively, the formation and photo-destruction rates of H<sub>2</sub>. In the case of the ISM, $`R`$ $``$ 3$`\times 10^{17}`$ s<sup>-1</sup> cm<sup>-3</sup> and $`\beta (0)`$ $`5\times 10^{10}`$ s<sup>-1</sup>. We use the ISM value of $`R`$ scaled by the dust content measured in the systems. It can be seen from Table 6 that for a moderate radiation field (like the ISM mean field) $`n_\mathrm{H}`$ derived using C ii for the CNM like parameters is usually higher than the upper limit obtained based on the H<sub>2</sub> equilibrium formation. In addition, the expected electron density for the $`T`$ = 100 K gas (assuming a $`n_\mathrm{e}/n_\mathrm{H}`$ ratio as seen in CNM) is higher than 10<sup>-2</sup> cm<sup>-3</sup>. At such electron densities, C i should be detectable. This is inconsistent with the non-detection of C i in these systems. This problem of CNM gas producing very small amount of C i is already recognized in the literature (Liszt et al. 2002; Wolfe et al. 2003).
We notice that the absence of C i and H<sub>2</sub> in these systems is consistent with the gas originating either from the WNM gas or from the ionized gas. As pointed out above, the strength of the Al iii absorption lines seen in these systemsare consistent with the gas density being less than the one expected for the CNM.
Thus, if one uses only C ii absorption line then the results are consistent with these systems originating from CNM gas. However the absence of H<sub>2</sub> and C i absorption lines together with the presence of Al iii following the profiles of singly ionized gas is inconsistent with standard CNM solutions. Thus most of the DLAs without H<sub>2</sub> are consistent with them originating from low density, high temperature and partially ionized gas.
### 6.5 $`z_{\mathrm{abs}}`$= 1.944 toward Q $`1157+014`$
Some of our conclusions of the previous Section can be ascertained in the case of the $`z_{\mathrm{abs}}`$= 1.944 system toward Q $`1157+014`$ as we have additional information on the kinetic temperature based on 21 cm absorption (Wolfe et al. 1981). The estimated spin temperature based on the recent measurement of $`N`$(H i) is $`T`$ = 865$`\pm `$190 K (Kanekar & Chengalur 2003). Fig. LABEL:vp1157 shows the velocity plot of selected absorption lines in this system. The distribution of neutral, singly ionized and doubly ionized species in velocity space can be visualized easily. Using a Gaussian profile, we also show, in the bottom panel, the velocity range over which 21 cm absorption is seen. Whereas the UV absorption lines spread over more than 100 km s<sup>-1</sup>, the 21 cm absorption originates from only a few of the components. There are components that do not possess 21cm absorption but show absorption due to UV transitions including that of C ii. As expected, C ii traces a wide range of physical conditions. Assuming the temperature of the gas that is producing the 21 cm absorption feature to be 100 K (or 200 K), we can derive from the 21 cm observation the H i column density in the component lying along the line of sight: $`7\times 10^{20}`$ cm<sup>-2</sup> (or 1.4$`\times 10^{21}`$ cm<sup>-2</sup>). This is approximately 15% (30%) of the total $`N`$(H i) measured from the damped Lyman-$`\alpha `$ line. This means that 85% (70%) of $`N`$(H i) along the line of sight is either warm or hot. From the Si ii$`\lambda `$1808 profile, we notice that $``$60% of $`N`$(Si ii) originate in the velocity space covered by the 21cm profile. We also notice that considerable fraction of Al iii absorption originate from the velocity range covered by 21 cm absorption. Thus warm and ionized components seem to be co-spacial with the cold gas responsible for the 21 cm absorption. In addition we also notice Al iii components with Si ii and C ii absorption well separated from the 21 cm component. Thus part of the C ii absorption seen here originate from the WNM or WIM. Therefore, the non-detection of H<sub>2</sub> in this system can be easily explained as a consequence of most of the gas being at high temperature (and hence low density). The absence of H<sub>2</sub> from the cold 21 cm absorbing component could just be due to the low $`N`$(H i) associated with this component and the relatively low dust depletion.
## 7 Discussion and conclusions
We have studied the physical conditions in damped Lyman-$`\alpha `$ systems (DLAs) using a sample of 33 systems toward 26 QSOs acquired for a recently completed survey of H<sub>2</sub> in DLAs by Ledoux et al. (2003). We use standard techniques to estimate the physical conditions prevailing in the gas. In this Section, we discuss some of the results and related issues.
### 7.1 High pressure of the H<sub>2</sub> gas
Our study shows that the H<sub>2</sub> components in DLAs trace Cold gas (153$`\pm 78`$ K) with relatively high pressure. The pressure in individual components (measured assuming a radiation field similar to our Galaxy) is in the range 824$``$30,000 cm<sup>-3</sup> K, a large fraction of the components being at high pressure. 42%, 20%, and 8% of the components have pressure in excess of 3000 cm<sup>-3</sup> K, 5000 cm<sup>-3</sup> K and 10<sup>4</sup> cm<sup>-3</sup> K, respectively. Based on the profiles of singly ionized species we note the H<sub>2</sub> components arise in gas with a wide range of molecular content and ionization state much like what we see in the Galactic ISM.
This is not unexpected. Indeed, in the framework of a galactic two-phase medium, the stable pressure range for the gas is 460$`\mathrm{P}/\mathrm{k}(\mathrm{cm}^2\mathrm{K})`$ 1750 (Wolfe et al. 2003). Clearly the pressure we derive in the H<sub>2</sub> components are much higher than this. From Table 3 of Wolfire et al. (1995) it can be seen that for a given metallicity an increase in the dust-to-gas ratio can lead to an increase in the allowed range of pressure, whereas an increase in the metallicity reduces the allowed range of pressure due to enhanced cooling. For conditions typical of DLAs, that is for metallicities of Z = $`1.0`$ and a dust-to-gas ratio ten times smaller than in the Galaxy, the stable pressure range is 1800-13000 cm<sup>-3</sup> s<sup>-1</sup> (Wolfire et al. 1995). Note that in the absence of any confining medium (or pressure equilibrium between different components) we expect such a high pressure gas to survive only for a short period of time (with a typical hydrodynamical time-scale of 10<sup>6</sup> years).
The pressure we infer depends very much on the intensity of the radiation field. A larger intensity implies and excess of UV pumping (on top of what we assume in our analysis) which, if taken into account, should reduce the hydrogen density derived using the C i fine-structure lines. At the same time, the temperature of the gas will increase due to photo-heating.
### 7.2 H<sub>2</sub> content
Ledoux et al. (2003) found that approximately 13$``$20% of DLAs show H<sub>2</sub> absorption lines with most of the H<sub>2</sub> components having column densities in the range $`16.0`$ log $`N`$(H<sub>2</sub>)(cm<sup>-2</sup>$`19.0`$. In the case of the Galactic ISM, only a minor fraction of the clouds fall in this range. This is expected because, above log $`N`$(H<sub>2</sub>)$``$ 16.0, self-shielding drastically decreases the photo-dissociation rate. On the contrary, in the LMC and SMC, a large fraction of lines of sight have $`16.0`$ log $`N`$(H<sub>2</sub>$`19.0`$ (see Tumlinson et al., 2002). Thus the trend noticed in DLAs could just be a generic feature of gas with low dust content and metallicity. It is however important to remember that the molecular fraction given in Ledoux et al. (2003) is an average over the whole line of sight. The actual molecular fraction in individual components may be much larger. Thus the low values of $`N`$(H<sub>2</sub>) that are observed could just be a consequence of low $`N`$(H i) in the corresponding individual components. Indeed, consistent models of DLAs (Srianand et al. 2005) require H i column densities much less than the total $`N`$(H i) measured in DLAs with H<sub>2</sub>. In addition, the absence of 21 cm absorption can be reconciled if the H<sub>2</sub> components have only part of the total H i (see below).
### 7.3 21 cm absorption
We have shown that detecting H<sub>2</sub> and C i absorption lines is an efficient way to trace the cold neutral gas in DLAs. H i 21 cm absorption line provides an independent way of detecting the CNM gas in DLAs. The detectability of 21 cm absorption line depends only on the amount of cold gas along the line of sight and the covering factor of the radio source. Thus, for a compact background source one can detect CNM gas with 21 cm without any bias from dust content or metallicity. At $`z_{\mathrm{abs}}`$$``$ 2, seven systems have been searched for 21 cm absorption and none have been detected (Kanekar & Chengalur 2003). Assuming $`T`$ = 200 K for the CNM gas, these authors estimated the filling factor of the CNM gas to be $`0.3`$. This is consistent with what we derive from our H<sub>2</sub> survey. Over the redshift range covered by our survey there are three cases for which information on 21 cm absorption and H<sub>2</sub> content are available. To our surprise there seems to be no correlation between 21 cm absorption and H<sub>2</sub> absorption.
The $`z_{\mathrm{abs}}`$= 2.811 system toward PKS 0528$``$255 show H<sub>2</sub> absorption in two distinct components without any corresponding 21 cm absorption (Carilli et al. 1996). The upper limit on $`\tau `$(21cm) gives $`N`$(H i)$`5\times 10^{20}`$ cm<sup>-2</sup> if the kinetic temperature is similar to $`T`$(OPR) we measure. Thus H<sub>2</sub> and 21 cm observations can be consistent with one another if no more than 20% of the total H i column density is associated with the H<sub>2</sub> component. The $`z_{\mathrm{abs}}`$= 1.944 system toward Q 1157$`+`$014 (Wolfe, Briggs & Jauncy, 1981; and discussion above) and $`z_{\mathrm{abs}}`$= 2.04 towards PKS 0458-020 (Briggs et al., 1989; Ge & Bechtold, 1999) show strong 21 cm absorption but no C i or H<sub>2</sub> absorption. In these systems the absence of H<sub>2</sub> and C i absorption could be either due to low density in the cold H i component or to the presence of higher ambient radiation field.
### 7.4 ii absorption
As pointed out before, C ii is detected in all the systems in which H<sub>2</sub> is seen. In fact, C ii is also detected in the three systems that show 21 cm absorption discussed in the previous Section. C ii being the dominant ion of Carbon in the neutral gas, it is natural to expect C ii associated with both 21 cm and H<sub>2</sub> absorption. However, C ii is readily detected in warm neutral gas and even in ionized gas. The interpretation of the origin of the C ii absorption is not as straightforward as in the case of H<sub>2</sub> and C i. Thus, the nature of systems that do not show 21 cm absorption and/or H<sub>2</sub> absorption is a matter of debate. The systems that show C ii in our sample are consistent with them originating from the CNM gas. However, the absence of C i and H<sub>2</sub> (if we take the depletion as an indicator of the presence of dust) and the presence of Al iii are also consistent with C ii absorption originating from the warm/partially ionized gas. Thus the frequency of occurrence of C ii provides a liberal upper limit on the CNM covering factor. A detailed investigation taking into account the constraints on the ionization state of the gas based on N ii, Fe iii or Al iii will be important to derive the exact covering factor of CNM gas.
### 7.5 Star-Formation Rate
One of the main driver for the study of DLAs is to find out a way to recover the global star-formation history in a typical, moderately star-forming environment. The importance of DLAs in the paradigm of hierarchical structure formation can be appreciated from the fact that the mass density of baryonic matter in DLAs at $`z_{\mathrm{abs}}3`$ is similar to that of stars at present epochs (Wolfe, 1995). Studies of Lyman-$`\alpha `$ and UV continuum emission from galaxies associated with DLAs usually result in star formation rates (or upper limits) of a few M yr<sup>-1</sup> (Fynbo et al., 1999; Bunker et al., 1999; Kulkarni et al., 2001).
Wolfe et al. (2003a, 2003b) have proposed a novel idea of using the C ii cooling rate to infer the SFR in DLA galaxies. The idea is that if one assumes thermal equilibrium then the cooling rate inferred from C ii should be equal to the heating rate (through UV photons, Cosmic rays etc.,) driven by the local star-formation activity. Their detailed study suggests the star formation rate density of DLAs at high-$`z`$ could be as high as that inferred based on Lyman break Galaxies. We confirm the presence of C ii in $``$ 50% of the DLAs in our sample. However, it is clear from the above discussion that one needs to unveil the nature of the partially ionized gas in order to have a handle on the heating rates.
In the local universe, star-formation is always related to molecular clouds. If DLAs are star-forming regions then the local star-formation has to be related to the mass of the molecular gas. Our survey shows that 13$``$20% of DLAs are associated with H<sub>2</sub> in absorption. We have not detected CO in any of these systems and HD is detected in only one system (Varshalavich et al., 2000). Clearly the dark molecular clouds where stars form in our Galaxy are not seen along QSO lines of sight. The UV radiation field inferred from the H<sub>2</sub> high-J excitation is similar to the Galactic mean field. Following Wolfe et al. (2003a,2003b) the SFR per unit comoving volume for DLAs is,
$`\dot{\rho }_{}`$ $`=`$ $`An_{co}(z)<\dot{\xi }(z)>`$ (9)
$`=`$ $`f_d<\dot{\xi }(z)>({\displaystyle \frac{A}{A_p}}){\displaystyle \frac{dN}{dX}}`$
where, $`<\dot{\xi }(z)>`$ is the average SFR per unit area at redshift $`z`$ and $`A,A_\mathrm{p}`$ and $`\frac{dN}{dX}`$ are average physical cross-sectional, respectively, area, average projected area and number density of absorbers per unit absorption distance interval. $`f_\mathrm{d}`$ is the fraction of DLAs in which the UV radiation field is similar to the Galactic UV background (i.e., 0.13$``$0.20). Here we use the fact that our H<sub>2</sub> sample is a randomly chosen sub-sample of the whole population of DLAs and the presence of H<sub>2</sub> is independent of $`N`$(H i).
For an Einstein-de Sitter cosmology, $`\frac{dN}{dX}=3\times 10^5`$ for the mean redshift of our sample (Storrie-Lombardi & Wolfe, 2000). Assuming $`H_0`$ = 75 km s<sup>-1</sup>Mpc<sup>-1</sup>, $`A/A_\mathrm{p}=2`$ and $`<\dot{\xi }(z)>=4\times 10^3`$ M yr<sup>-1</sup> kpc<sup>-2</sup> (typical for our Galaxy, see Kennicutt, 1998) we derive $`\dot{\rho }_{}0.03`$ at $`z_{\mathrm{abs}}`$= 2.5. This crude estimate already gives half the star formation rate density measured in Lyman break galaxies (Steidel et al., 1999). Recently, Hirashita & Ferrera (2005) have also arrived at similar conclusion. Thus it is of the utmost importance to understand the physics of the ISM in high-$`z`$ DLAs in order to derive the cosmic star-formation budget correctly.
## acknowledgements
Results presented in this work are based on observations carried out at the European Southern Observatory (ESO) under prog. ID No. 65.P-0038, 65.O-0063, 66.A-0624, 67.A-0078, 68.A-0600 68.A-0106 and 70.A-0017 with the UVES spectrograph installed on the Very Large Telescope (VLT) at Cerro Paranal Observatory in Chile. RS and PPJ gratefully acknowledge support from the Indo-French Centre for the Promotion of Advanced Research (Centre Franco-Indien pour la Promotion de la Recherche Avancรฉe) under contract No. 3004-3. GJF acknowledges the support of the NSF through AST 00-71180 and NASA with grant NAG5-12020. GJF and RS acknowledge the support from the DST/INT/US(NSF-RP0-115)/2002. GS would like to thank CCS, University of Kentucky for their two years of support. The hospitality of IUCAA is gratefully acknowledged.
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# Absence of slow transients, and the effect of imperfect vertical alignment, in turbulent Rayleigh-Bรฉnard convection
## 1 Introduction
Turbulent convection in a fluid heated from below, known as Rayleigh-Bรฉnard convection (RBC), has been under intense study for some time \[for reviews, see e.g. \[Siggia(1994), Kadanoff(2001), Ahlers, Grossmann & Lohse (2002)\]\]. A central prediction of models for this system \[\[Kraichnan(1962), Castaing et al.(1998), Shraiman and Siggia (1990), Grossmann & Lohse (2001)\]\] is the heat transported by the fluid. It is usually described in terms of the Nusselt number
$$๐ฉ=\frac{QL}{A\lambda \mathrm{\Delta }T}$$
(1)
where $`Q`$ is the heat current, $`L`$ the cell height, $`A`$ the cross-sectional area, $`\lambda `$ the thermal conductivity, and $`\mathrm{\Delta }T`$ the applied temperature difference. The Nusselt number depends on the Rayleigh number
$$R=\alpha g\mathrm{\Delta }TL^3/\kappa \nu $$
(2)
and on the Prandtl number
$$\sigma =\nu /\kappa .$$
(3)
Here $`\alpha `$ is the isobaric thermal expansion coefficient, $`g`$ the acceleration of gravity, $`\kappa `$ the thermal diffusivity, and $`\nu `$ the kinematic viscosity.
An important feature of turbulent RBC is the existence of a large-scale circulation (LSC) of the fluid \[\[Krishnamurty & Howard(1981)\]\]. For cylindrical samples of aspect ratio $`\mathrm{\Gamma }L/D1`$ the LSC is known to consist of a single cell, with fluid rising along the wall at some azimuthal location $`\theta `$ and descending along the wall at a location $`\theta +\pi `$ \[see, for instance, \[Qiu and Tong (2001a)\]\]. As $`\mathrm{\Gamma }`$ decreases, the nature of the LSC is believed to change. For $`\mathrm{\Gamma }\stackrel{<}{_{}}0.5`$ it is expected \[\[Verzicco & Camussi(2003), Stringano & Verzicco(2005), Sun et al. (2005)\]\] that the LSC consists of two or more convection cells, situated vertically one above the other. Regardless of the LSC structure, the heat transport in turbulent RBC is mediated by the emission of hot (cold) volumes of fluid known as โplumesโ from a more or less quiescent boundary layer above (below) the bottom (top) plate. These plumes are swept away laterally by the LSC and rise (fall) primarily near the side wall. Their buoyancy helps to sustain the LSC.
In a recent paper \[Chillร et al.(2004)\] reported measurements using a cylindrical sample of water with $`\sigma 2.33`$ and with $`L=1`$ m and $`D=0.5`$ m for $`R10^{12}`$. Their sample thus had an aspect ratio $`\mathrm{\Gamma }0.5`$ at the borderline between a single-cell and a multi-cell LSC. They found exceptionally long relaxation times of $`๐ฉ`$ that they attributed to a switching of the LSC structure between two states. Multi-stability was observed also in Nusselt-number measurements by \[Roche et al. (2004)\] for a $`\mathrm{\Gamma }=0.5`$ sample (see also \[Nikolaenko et al. (2005)\] for a discussion of these data). Chillรก et al. also found that $`๐ฉ`$ was reduced by tilting the sample through an angle $`\beta `$ relative to gravity by an amount given approximately by $`๐ฉ(\beta )/๐ฉ(0)12\beta `$ when $`\beta `$ is measured in radian. A reduction by two to five percent of $`๐ฉ`$ (depending on $`R`$) due to a tilt by $`\beta 0.035`$ of a $`\mathrm{\Gamma }=0.5`$ sample was reported as well recently by \[Sun et al. (2005)\], although in that paper the $`\beta `$-dependence of this effect was not reported. Chillรก et al. developed a simple model that yielded a depression of $`๐ฉ`$ for the two-cell structure that was consistent in size with their measurements. Their model also assumes that no depression of $`๐ฉ`$ should be found for a sample of aspect ratio near unity where the LSC is believed to consist of a single convection cell; they found some evidence to support this in the work of \[Belmonte et al. (1995)\]. Indeed, recent measurements by \[Nikolaenko et al. (2005)\] for $`\mathrm{\Gamma }=1`$ gave the same $`๐ฉ`$ within 0.1 percent for a level sample and a sample tilted by 0.035 rad.
In this paper we report on a long-term study of RBC in a cylindrical sample with $`\mathrm{\Gamma }1`$. As expected, we found no long relaxation times because the LSC is uniquely defined. The establishment of a statistically stationary state after a large change of $`R`$ occurred remarkably quickly, within a couple of hours, and thereafter there were no further long-term drifts over periods of many days.
We also studied the orientation $`\theta _0`$ of the circulation plane of the LSC by measuring the side-wall temperature at eight azimuthal locations \[\[Brown et al. (2005b)\]\]. With the sample carefully leveled (i.e. $`\beta =0`$) we found $`\theta _0`$ to change erratically, with large fluctuations. There were occasional relatively rapid reorientations, as observed before by \[Sreenivasan et al. (2002)\]. The reorientations usually consisted of relatively rapid rotations, and rarely were reversals involving the cessation of the LSC followed by its re-establishment with a new orientation. This LSC dynamics yielded a broad probability distribution-function $`P(\theta _0)`$, although a preferred orientation prevailed. When the sample was tilted relative to gravity through an angle $`\beta `$, a well defined new orientation of the LSC circulation plane was established, $`P(\theta _0)`$ became much more narrow, and virtually all meandering and reorientation of the LSC was suppressed.
We found that $`๐ฉ`$ was reduced very slightly by tilting the sample. We obtained $`๐ฉ(\beta )=๐ฉ_0[1(3.1\pm 0.1)\times 10^2|\beta |]`$. This effect is about a factor of 50 smaller than the one observed by Chillรก et al. for their $`\mathrm{\Gamma }=0.5`$ sample.
From side-wall-temperature measurements at two opposite locations we determined time cross-correlation functions $`C_{i,j}`$. The $`C_{i,j}`$ had a peak that could be fitted well by a Gaussian function, centered about a characteristic time $`t_1^{cc}`$ that we interpreted as corresponding to the transit time needed by long-lived thermal disturbances to travel with the LSC from one side of the sample to the other, i.e. to half a turnover time of the LSC. We found that the $`\beta `$-dependence of the corresponding Reynolds number $`R_e^{cc}`$ is given by $`R_e^{cc}(\beta )=R_e^{cc}(0)\times [1+(1.85\pm 0.21)|\beta |(5.9\pm 1.7)\beta ^2]`$. A similar result was obtained from the auto-correlation functions of individual thermometers. Thus there is an $`๐ช(1)`$ effect of $`\beta `$ on $`R_e`$, and yet the effect of $`\beta `$ on $`๐ฉ`$ was seen to be nearly two orders of magnitude smaller. We also determined the temperature amplitude $`\delta `$ of the azimuthal temperature variation at the mid-plane. We expect $`\delta `$ to be a monotonically increasing function of the speed of the LSC passing the mid-plane, i.e. of the Reynolds number. We found $`\delta (\beta )=\delta (0)\times [1+(1.84\pm 0.45)|\beta |(3.1\pm 3.9)\beta ^2]`$. Thus, for small $`\beta `$ its $`\beta `$-dependence is very similar to that of the Reynolds number.
From the large effect of $`\beta `$ on $`R_e`$ and the very small effect on $`๐ฉ`$ we come to the important conclusion that the heat transport in this system is not influenced significantly by the strength of the LSC. This heat transport thus must be determined primarily by the efficiency of instability mechanisms in the boundary layers. It seems reasonable that these mechanisms should be nearly independent of $`\beta `$ when $`\beta `$ is small. This result is consistent with prior measurements by \[Ciliberto et al.(1997)\], who studied the LSC and the Nusselt number in a sample with a rectangular cross section. They inserted vertical grids above (below) the bottom (top) plate that suppressed the LSC, and found that within their resolution of a percent or so the heat transport was unaltered. Their shadowgraph visualizations beautifully illustrate that the plumes are swept along laterally by the LSC when there are no grids and rise or fall vertically due to their buoyancy in the presence of the grids. \[Ciliberto et al.(1997)\] also studied the effect of tilting their rectangular sample by an angle of 0.17 rad. Consistent with the very small effect of tilting on $`๐ฉ`$ found by us, they found that within their resolution the heat transport remained unaltered.
We observed that the sudden reorientations of the LSC that are characteristic of the level sample are strongly suppressed by even a small tilt angle.
## 2 Apparatus and Data Analysis
For the present work we used the โlargeโ and the โmediumโ sample and apparatus described in detail by \[Brown et al. (2005a)\]. Copper top and bottom plates each contained five thermistors close to the copper-fluid interface. The bottom plate had imbedded in it a resistive heater capable of delivering up to 1.5 kW uniformly distributed over the plate. The top plate was cooled via temperature-controlled water circulating in a double-spiral channel. For the Nusselt-number measurements a temperature set-point for a digital feedback regulator was specified. The regulator read one of the bottom-plate thermometers at time intervals of a few seconds and provided appropriate power to the heater. The top-plate temperature was determined by the temperature-controlled cooling water from two Neslab RTE740 refrigerated circulators.
Each apparatus was mounted on a base plate that in turn was supported by three legs consisting of long threaded rods passing vertically through the plate. The entire apparatus thus could be tilted by an angle $`\beta `$ relative to the gravitational acceleration by turning one of the rods. The maximum tilt angle attainable was 0.12 (0.21) rad for the large (medium) sample.
The Nusselt number was calculated using the temperatures recorded in each plate and the power dissipated in the bottom-plate heater. The side wall was plexiglas of thickness 0.64 cm (0.32 cm) for the large (medium) sample. It determined the length $`L`$ of the sample. Around a circumference the height was uniformly $`50.62\pm 0.01`$ cm ( $`24.76\pm 0.01`$ cm) for the large (medium) sample. The inside diameter was $`D=49.70\pm 0.01`$ cm ( $`D=24.84\pm 0.01`$ cm) for the large (medium) sample. The end plates had anvils that protruded into the side wall, thus guaranteeing a circular cross section near the ends. For the large sample we made measurements of the outside diameter near the half-height after many months of measurements and found that this diameter varied around the circumference by less than 0.1%.
Imbedded in the side wall and within 0.06 cm of the fluid-plexiglas interface were eight thermistors, equally spaced azimuthally and positioned vertically at half height of the sample. They yielded a relatively high (low) temperature reading at the angular positions where there was up-flow (down-flow) of the LSC. A fit of
$$T_i=T_c+\delta cos(i\pi /4\theta _0),i=0,\mathrm{},7$$
(4)
yielded the mean center temperature $`T_c`$, the angular orientation $`\theta _0`$ of the LSC (relative to the location of thermistor 0), and a measure $`\delta `$ of the LSC strength.
We expect the size of $`\delta `$ to be determined by the heat transport across a viscous boundary layer separating the LSC from the side wall. Thus $`\delta `$ should be a monotonically increasing function of the LSC Reynolds number $`R_e`$ because the boundary-layer thickness is expected to decrease with $`R_e`$ as $`1/R_e^{1/2}`$, and because the azimuthal temperature variation carried by the LSC near the boundary layer increases with $`R`$ and thus with $`R_e`$. However, the precise relationship between $`\delta `$ and $`R_e`$ is not obvious. Experimentally we find, over the range $`5\times 10^9<R<10^{11}`$ and for the large sample, that $`\delta `$ is related to $`R`$ by an effective power law $`\delta R^{0.81}`$, whereas $`R_eR^{0.50}`$ in this range, yielding $`\delta R_e^{1.62}`$. We would then expect that $`\delta `$ and $`R_e`$ will have a similar dependence on $`\beta `$ (at least for small $`\beta `$), albeit possibly with somewhat different coefficients.
From time series of the $`T_i(t)`$ taken at intervals of a few seconds and covering at least one day we determined the cross-correlation functions $`C^{i,j}(\tau )`$ corresponding to signals at azimuthal positions displaced around the circle by $`\pi `$ (i.e. $`j=i+4`$). These functions are given by
$$C^{i,j}(\tau )=[T_i(t)T_i(t)_t]\times [T_j(t+\tau )T_j(t)_t]_t.$$
(5)
We also calculated the auto-correlation functions corresponding to $`i=j`$ in Eq. 5, for all eight thermometers.
Initially each sample was carefully leveled so that the tilt angle relative to gravity was less than $`10^3`$ radian. Later it was tilted deliberately to study the influence of a non-zero $`\beta `$ on the heat transport.
The fluid was water at 40C where $`\alpha =3.88\times 10^4`$ K<sup>-1</sup>, $`\kappa =1.53\times 10^3`$ cm<sup>2</sup>/s, and $`\nu =6.69\times 10^3`$ cm<sup>2</sup>/s, yielding $`\sigma =4.38`$.
## 3 The Nusselt number of a vertical sample
### 3.1 Initial transients
In Fig. 1a we show the initial evolution of the top and bottom temperatures of the large sample in a typical experiment. Initially the heat current was near zero and $`T_b`$ and $`T_t`$ were close to 40C. The sample had been equilibrated under these conditions for over one day. Near $`t=0.6`$ h a new temperature set point of 50C was specified for the bottom plate, and the circulator for the top plate was set to provide $`T_t30^{}`$C. From Fig. 1a one sees that there were transients that lasted until about 0.9 h (1.2 h) for $`T_b`$ ($`T_t`$). These transients are determined by the response time and power capability of the bottom-plate heater and the top-plate cooling water and are unrelated to hydrodynamic phenomena in the liquid. Figures 1b and c show the evolution of the heat current. After the initial rapid rise until $`t0.8`$ h the current slowly evolved further to a statistically stationary value until $`t3`$ h. A fit of the exponential function $`Q(t)=Q_{\mathrm{}}\mathrm{\Delta }Qexp(t/\tau _Q)`$ to the data for $`t>1.2`$ h is shown by the solid line in Fig. 1c and yielded a relaxation time $`\tau _Q=0.48\pm 0.04`$ h. We attribute this transient to the evolution of the fluid flow. It is interesting to compare $`\tau _Q`$ with intrinsic time scales of the system. The vertical thermal diffusion time $`\tau _vL^2/\kappa `$ is 467 hours. Obviously it does not control the establishment of the stationary state. If we consider that it may be reduced by a factor of $`1/๐ฉ`$ with $`๐ฉ=263`$, we still obtain a time sale of 1.78 hours that is longer than $`\tau _Q`$. We believe that the relatively rapid equilibration is associated with the establishment of the top and bottom boundary layers that involve much shorter lengths $`l_t`$ and $`l_b`$. It also is necessary for the large-scale circulation to establish itself; but, as we shall see, its precise Reynolds number is unimportant for the heat transport. In addition, the LSC can be created relatively fast since this is not a diffusive process.
### 3.2 Results under statistically stationary conditions
Figure 1 shows the behavior of the system only during the first six hours and does not exclude the slow transients reported by Chillร et al. that occurred over time periods of $`๐ช(10^2)`$ hours. Thus we show in Fig. 2 results for $`๐ฉ/R^{1/3}`$ from a run using the large sample that was continued under constant externally imposed conditions for nine days. Each point corresponds to a value of $`๐ฉ`$ based on a time average over two hours of the plate temperatures and the heat current. Note that the vertical range of the entire graph is only 0.8 %. Thus, within a small fraction of 1 %, the results are time independent. Indeed, during nearly a year of data acquisition for a $`\mathrm{\Gamma }=1`$ sample at various Rayleigh numbers, involving individual runs lasting from one to many days, we have never experienced long-term drifts or changes of $`๐ฉ`$ after the first few hours. This differs dramatically from the observations of Chillร et al. who found changes by about 2 % over about 4 days. We conclude that the slow transients observed by them for their $`\mathrm{\Gamma }=0.5`$ sample do not occur for $`\mathrm{\Gamma }1`$.
To document further the stationary nature of the system, we compared results from the large sample for $`๐ฉ`$ obtained from many runs, each of one to ten daysโ duration, over a period of about five months \[\[Nikolaenko et al. (2005), Funfschilling et al. (2005)\]\]. The scatter of the data at a given $`R`$ is only about 0.1%. This excellent reproducibility would not be expected if there were slow transients due to transitions between different states of the LSC.
Although work in our laboratory with other aspect ratios has been less extensive, we also have not seen any evidence of drifts or transients for the larger $`\mathrm{\Gamma }=1.5,2,3,`$ and 6 \[\[Funfschilling et al. (2005), Brown et al. (2005a)\]\] nor for the smaller $`\mathrm{\Gamma }=0.67,0.43,`$ and 0.28 \[\[Nikolaenko et al. (2005), Brown et al. (2005a)\]\]. It may be that $`\mathrm{\Gamma }=0.5`$, being near the borderline between a single-cell LSC and more complicated LSC structures \[\[Verzicco & Camussi(2003), Stringano & Verzicco(2005), Sun et al. (2005)\]\], is unique in this respect.
In Fig. 3 we compare results for $`๐ฉ/R^{1/3}`$ from our large sample \[\[Nikolaenko et al. (2005)\]\] with those reported by Chillร et al. (stars). Our results are larger by about 15%. To find a reason for this difference, we first look at the $`\mathrm{\Gamma }`$ and $`\sigma `$ dependence. The open (solid) circles represent our data for $`\sigma =4.38`$ and $`\mathrm{\Gamma }=0.67(0.43)`$ and show that the dependence of $`๐ฉ`$ on $`\mathrm{\Gamma }`$ is not very strong. The open squares (diamonds) are our results for $`\mathrm{\Gamma }=0.67`$ and $`\sigma =5.42(3.62)`$ and indicate that $`๐ฉ`$ actually increases slightly with $`\sigma `$. Thus the lower values of $`๐ฉ`$ (compared to ours) obtained by Chillรก et al. for $`\sigma =2.3`$ and $`\mathrm{\Gamma }=0.5`$ can not be explained in terms of the $`\mathrm{\Gamma }`$ and $`\sigma `$ dependence of $`๐ฉ`$. Some of the difference can be attributed to non-Boussinesq effects that tend to reduce $`๐ฉ`$ \[\[Funfschilling et al. (2005)\]\]. However, for the largest $`\mathrm{\Delta }T`$ used by Chillรก et al. (31C) we expect this effect to be somewhat less than 1 % \[\[Funfschilling et al. (2005)\]\]. Finally, the effect of the finite conductivity of the top and bottom plates comes to mind. This can reduce $`๐ฉ`$ by several % when $`\mathrm{\Delta }T`$ is large \[\[Chaumat et al.(2002), Verzicco(2004), Brown et al. (2005a)\]\], but it is difficult to say precisely by how much. It seems unlikely that this effect can explain the entire difference, particularly at the smaller $`R`$ (and thus $`\mathrm{\Delta }T`$) where it is relatively small.
## 4 Tilt-angle dependence of the Nusselt number
In Fig. 4 we show results for $`๐ฉ`$ from the large sample at $`R=9.43\times 10^{10}`$. Each data point was obtained from a two-hour average of measurements of the various temperatures and of $`Q`$. Three data sets, taken in temporal succession, for tilt angles $`\beta =0.000,0.087`$, and 0.122 are shown. All data were normalized by the mean of the results for $`\beta =0`$. Typically, the standard deviation from the mean of the data at a given $`\beta `$ was 0.13%. The vertical dotted lines and the change in the data symbols show where $`\beta `$ was changed. One sees that tilting the cell caused a small but measurable reduction on $`๐ฉ`$. In Fig. 5 we show the mean value for each tilt angle, obtained from runs of at least a dayโs duration a each $`\beta `$, as a function of $`|\beta |`$. One sees that $`๐ฉ`$ decreases linearly with $`\beta `$. A fit of a straight line to the data yielded
$$๐ฉ(\beta )=๐ฉ_0[1(3.1\pm 0.1)\times 10^2|\beta |].$$
(6)
with $`๐ฉ_0=273.5`$. Simlar results for the medium cell are compared with the large-cell results in Fig. 6. At the smaller Rayleigh number of the medium sample the effect of $`\beta `$ on $`๐ฉ`$ is somewhat less. Because the effect of $`\beta `$ on $`๐ฉ`$ is so small, we did not make a more detailed investigation of its Rayleigh-number dependence.
Chillร et al. proposed a model that predicts a significant tilt-angle effect on $`๐ฉ`$ for $`\mathrm{\Gamma }=0.5`$ where they assume the existence of two LSC cells, one above the other. They also assumed that there would be no effect for $`\mathrm{\Gamma }=1`$ where there is only one LSC cell. Although we found an effect for our $`\mathrm{\Gamma }=1`$ sample, we note that it is a factor of about 50 smaller than the effect observed by Chillร et al. for $`\mathrm{\Gamma }=0.5`$.
## 5 Tilt-angle dependence of the large-scale circulation
### 5.1 The orientation
In Fig. 7 we show the angular orientation $`\theta _0`$ (a) and the temperature amplitude $`\delta `$ (b) of the LSC. For the first 8000 sec shown in the figure, the sample was level ($`\beta =0.000\pm 0.001`$). One sees that $`\theta _0`$ varied irregularly with time. The probability-distribution function $`P(\theta _0)`$ is shown in Fig. 8 as solid dots. Essentially all angles are sampled by the flow, but there is a preferred direction close to $`\theta _0/2\pi =0.6`$. At t = 8000 sec, the sample was tilted through an angle $`\beta =0.087`$ radian. The direction of the tilt was chosen deliberately so as to oppose the previously prevailing preferred orientation. As a consequence one sees a sharp transition with a change of $`\theta _0`$ by approximately $`\pi `$. The temperature amplitude $`\delta `$ on average increased slightly, and certainly remained non-zero. From this we conclude that the transition took place via rotation of the LSC, and not by cessation that would have involved a reduction of $`\delta `$ to zero \[see \[Brown et al. (2005b)\]\]. We note that $`\theta _0(t)`$ fluctuated much less after the tilt. The results for $`P(\theta _0)`$ after the tilt are shown in Fig. 8 as open circles. They confirm that the maximum was shifted close to $`\theta _0=0`$, and that the distribution was much more narrow.
In Fig. 9 we show $`P(\theta _0)`$ for $`\beta =0.122,`$ (solid circles), 0.044 (open circles), and 0.026 (solid squares). One sees that a reduction of $`\beta `$ leads to a broadening of $`P(\theta _0)`$. The square root of the variance of data like those in Fig. 9 is shown in Fig. 10 on a logarithmic scale as a function of $`\beta `$ on a linear scale. Even a rather small tilt angle caused severe narrowing of $`P(\theta _0)`$.
### 5.2 The temperature amplitude
In Fig. 11 we show the temperature-amplitude $`\delta (\beta )`$ of the LSC as a function of $`\beta `$. As was the case for $`๐ฉ`$, the data are averages over the duration of a run at a given $`\beta `$ (typically a day or two). The solid (open) circles are for positive (negative) $`\beta `$. The data can be represented well by either a linear or a quadratic equation. A least-squares fit yielded
$$\delta (\beta )=\delta (0)\times [1+(1.84\pm 0.45)|\beta |(3.1\pm 3.9)\beta ^2]$$
(7)
with $`\delta (0)=0.164`$ K.
### 5.3 The Reynolds numbers
Using Eq. 5, we calculated the auto-correlation functions (AC) $`C^{i,i},i=0,\mathrm{},7`$, as well as the cross-correlation functions (CC) $`C^{i,j},j=(i+4)\%8,i=0,\mathrm{},7`$, of the temperatures measured on opposite sides of the sample. Typical examples are shown in Fig. 12. The CC has a characteristic peak that we associate with the passage of relatively hot or cold volumes of fluid at the thermometer locations. Such temperature cross-correlations have been shown, e.g. by \[Qiu and Tong (2001b)\] and \[Qiu and Tong (2002)\], to yield delay times equal to those of velocity-correlation measurements, indicating that warm or cold fluid volumes travel with the LSC. The function
$$C^{i,j}(\tau )=b_0exp\left(\frac{\tau }{\tau _0^{i,j}}\right)b_1exp\left[\left(\frac{\tau t_1^{i,j}}{\tau _1^{i,j}}\right)^2\right],$$
(8)
consisting of an exponentially decaying background (that we associate with the random time evolution of $`\theta _0`$) and a Gaussian peak, was fitted to the data for the CC. The fitted function is shown in Fig. 12 as a solid line over the range of $`\tau `$ used in the fit. It is an excellent representation of the data and yields the half turnover time $`๐ฏ/2=t_1^{i,j}`$ of the LSC. Similarly, we fitted the function
$$C^{i,i}(\tau )=b_0exp\left(\frac{\tau }{\tau _0^{i,i}}\right)+b_1exp\left[\left(\frac{\tau }{\tau _1^{i,i}}\right)^2\right]+b_2exp\left[\left(\frac{\tau t_2^{i,i}}{\tau _2^{i,i}}\right)^2\right]$$
(9)
to the AC data. It consists of two Gaussian peaks, one centered at $`\tau =0`$ and the other at $`\tau =t_2^{i,i}`$, and the exponential background. We interpret the location $`t_2^{i,i}`$ of the second Gaussian peak to correspond to a complete turnover time $`๐ฏ`$ of the LSC.
In terms of the averages $`<t_1^{i,j}>`$ and $`<t_2^{i,i}>`$ over all 8 thermometers or thermometer-pair combinations we define \[\[Qiu and Tong (2002), Grossmann & Lohse (2002)\]\] the Reynolds numbers
$$R_e^{cc}=(L/<t_1^{i,j}>)(L/\nu )$$
(10)
and
$$R_e^{ac}=(2L/<t_2^{i,i}>)(L/\nu ).$$
(11)
Here the length scale $`2L`$ was used to convert the turnover time $`๐ฏ`$ into a LSC speed $`2L/๐ฏ`$. For $`\mathrm{\Gamma }=1`$, the length $`4L`$ might have been used instead, as was done for instance by \[Lam et al.(2002)\]. This would have led to a Reynolds number larger by a factor of two. In Fig. 13a and b we show $`R_e^{cc}(|\beta |)`$ and $`R_e^{ac}(|\beta |)`$ respectively. The solid circles are for positive and the open ones for negative $`\beta `$. One sees that $`R_e^{cc}(|\beta |)`$ and $`R_e^{ac}(|\beta |)`$ initially grow linearly with $`\beta `$, but the data also reveal some curvature as $`|\beta |`$ becomes larger. Thus we fitted quadratic equations to the data and obtained
$$R_e^{cc}(\beta )=R_e^{cc}(0)\times [1+(1.85\pm 0.21)|\beta |(5.9\pm 1.7)\beta ^2]$$
(12)
and
$$R_e^{ac}(\beta )=R_e^{ac}(0)\times [1+(1.72\pm 0.38)|\beta |(4.1\pm 3.2)\beta ^2]$$
(13)
with $`R_e^{cc}(0)=10467\pm 43`$ and $`R_e^{ac}(0)=10565\pm 82`$ (all parameter errors are 67% confidence limits). The results for $`R_e^{cc}(0)`$ and $`R_e^{ac}(0)`$ are about 10% higher than the prediction by \[Grossmann & Lohse (2002)\] for our $`\sigma `$ and $`R`$. The excellent agreement between $`R_e^{cc}`$ and $`R_e^{ac}`$ is consistent with the idea that the CC yields $`๐ฏ/2`$ and that the AC gives $`๐ฏ`$. As expected (see Sect. 2), the $`\beta `$-dependences of both Reynolds numbers are the same within their uncertainties. It is interesting to see that the coefficients of the linear term also agree with the corresponding coefficient for $`\delta `$ (Eq. 7). This suggests that there may be a closer relationship between $`\delta `$ and $`R_e`$ than we would have expected a priori. However, the coefficient of the linear term in Eq. 12 or 13 is larger by a factor of about 50 than the corresponding coefficient for the Nusselt number in Eq. 6.
In Fig. 14 we show measurements of $`R_e^{cc}`$ and of $`\delta `$, each normalized by its value at $`\beta =0`$, as a function of $`\beta `$ for the medium sample and $`R=1.13\times 10^{10}`$. For this sample we were able to attain larger values of $`\beta `$ than for the large one. One sees that $`\delta `$ and $`R_e^{cc}`$ have about the same $`\beta `$ dependence for small $`\beta `$, but that $`\delta `$ then increases more rapidly than $`R_e^{cc}`$ as $`\beta `$ becomes large. Although we do not know the reason for this behavior, it suggests that the larger speed of the LSC enhances the thermal contact between the side wall and the fluid interior.
The Rayleigh-number dependence of $`R_e^{cc}`$ at constant $`\beta `$ is shown in Fig. 15. Here the open (solid) circles are from the medium (small) sample. There is consistency between the two samples, and the data can be described by a power law with a small negative exponent. The solid line is drawn to correspond to an exponent of $`1/6`$.
### 5.4 A model for the enhancement of the Reynolds number
As seen in Fig. 8, the LSC assumes an orientation for which gravity enhances the velocity above (below) the bottom (top plate), i.e. the LSC flows โuphillโ at the bottom where it is relatively warm and โdownhillโ at the top where it is relatively cold. This leads to an enhancement of the Reynolds number of the LSC. As suggested by \[Chillร et al.(2004)\], one can model this effect by considering the buoyancy force per unit area parallel to the plates. This force can be estimated to be $`\rho lg\beta \alpha \mathrm{\Delta }T/2`$ where $`l`$ is the boundary-layer thickness. It is opposed by the increase of the viscous shear stress across the boundary layer that may be represented by $`\rho \nu u^{}/l`$ where $`u^{}`$ is the extra speed gained by the LSC due to the tilt. Equating the two, substituting
$$l=L/(2๐ฉ),$$
(14)
solving for $`u^{}`$, using Eq. 2 for $`R`$, and defining $`R_e^{}(L/\nu )u^{}`$ one obtains
$$R_e^{}=\frac{R\beta }{8\sigma ๐ฉ^2}$$
(15)
for the enhancement of the Reynolds number of the LSC. From our measurements at large $`R`$ we found that $`R_e`$ \[\[Ahlers et al. (2005)\]\] and $`๐ฉ`$ \[\[Nikolaenko et al. (2005)\]\] can be represented within experimental uncertainty by
$`R_e`$ $`=`$ $`0.0345R^{1/2},`$ (16)
$`๐ฉ`$ $`=`$ $`0.0602R^{1/3},`$ (17)
giving
$$\frac{R_e^{}}{R_e}=1.00\times 10^3R^{1/6}\sigma ^1\beta .$$
(18)
For our $`\sigma =4.38`$ and $`R=9.43\times 10^{10}`$ one finds $`R_e^{}/R_e=3.4\beta `$, compared to the experimental value $`(1.9\pm 0.2)\beta `$ from $`R_e^{cc}`$ \[Eq. 12\] and $`(1.7\pm 0.4)\beta `$ from $`R_e^{ac}`$ \[Eq. 13\]. We note that the coefficient $`1.00\times 10^3`$ in Eq. 18 depends on the definition of $`R_e`$ given in Eqs. 10 an 11 that was used in deriving the result Eq. 16. If the length scale $`4L`$ had been used instead of $`2L`$ to define the speed of the LSC, as was done for instance by \[Lam et al.(2002)\], this coefficient would have been smaller by a factor of two, yielding near-perfect agreement with the measurements. In Fig. 15 one sees that the predicted dependence on $`R^{1/6}`$ also is in excellent agreement with the experimental results. However, such good agreement may be somewhat fortuitous, considering the approximations that were made in the model. Particularly the use of Eq. 14 for the boundary-layer thickness is called into question at a quantitative level by measurements of \[Lui & Xia(1998)\] that revealed a significant variation of $`l`$ with lateral position. In addition, it is not obvious that the thermal boundary-layer thickness $`l`$ should be used, as suggested by \[Chillร et al.(2004)\], to estimate the shear stress; perhaps the thickness of the viscous BL would be more appropriate.
In discussing their $`\mathrm{\Gamma }=0.5`$ sample, \[Chillร et al.(2004)\] took the additional step of assuming that the relative change due to a finite $`\beta `$ of $`๐ฉ`$ is equal to the relative change of $`R_e`$. For our sample with $`\mathrm{\Gamma }=1`$ this assumption does not hold. As we saw above, the relative change of $`๐ฉ`$ is a factor of about 50 less than the relative change of $`R_e`$. The origin of the (small) depression of the Nusselt number is not so obvious. Naively one might replace $`g`$ in the definition of the Rayleigh number by $`gcos(\beta )`$; but this would lead to a correction of order $`\beta ^2`$ whereas the experiment shows that the correction is of order $`\beta `$, albeit with a coefficient that is smaller than of order one. The linear dependence suggests that the effect of $`\beta `$ on $`๐ฉ`$ may be provoked by the change of $`R_e`$ with $`\beta `$, but not in a direct causal relationship.
## 6 Tilt-angle dependence of reorientations of the large-scale circulation
It is known from direct numerical simulation \[\[Hansen et al.(1991)\]\] and from several experiments \[\[Cioni et al.(1997), Niemela et al. (2001), Sreenivasan et al. (2002), Brown et al. (2005b)\]\] that the LSC can undergo relatively sudden reorientations. Not unexpectedly, we find that the tilt angle strongly influences the frequency of such events. For a level sample ($`\beta =0`$) we demonstrated elsewhere \[\[Brown et al. (2005b)\]\] that reorientations can involve changes of the orientation of the plane of circulation of the LSC through any angular increment $`\mathrm{\Delta }\theta `$, with the probability $`P(\mathrm{\Delta }\theta )`$ increasing with decreasing $`\mathrm{\Delta }\theta `$. Thus, in order to define a โreorientationโ, we established certain criteria. We required that the magnitude of the net angular change $`|\mathrm{\Delta }\theta |`$ had to be greater than $`\mathrm{\Delta }\theta _{min}=(2\pi )/8`$. In addition we specified that the magnitude of the net average azimuthal rotation rate $`|\dot{\theta }||\mathrm{\Delta }\theta /\mathrm{\Delta }t|`$ had to be greater than $`\dot{\theta }_{min}=0.1/๐ฏ`$ where $`๐ฏ`$ is the LSC turnover time and $`\mathrm{\Delta }t`$ is the duration of the reorientation (we refer to \[Brown et al. (2005b)\] for further details). Using these criteria, we found that the number of reorientation events $`n(\beta )`$ at constant $`R=9.43\times 10^{10}`$ decreased rapidly with increasing $`|\beta |`$. These results are shown in Fig. 16. It is worth noting that nearly all of these events are rotations of the LSC and very few involved a cessation of the circulation. A least-squares fit of the Gaussian function
$$n(\beta )=N_0exp[(\beta \beta _0)^2/w^2]$$
(19)
to the data yielded $`N_0=1.23\pm 0.06`$ events per hour, $`\beta _0=0.0093\pm 0.0010`$ rad, and $`w=0.0251\pm 0.0015`$ rad. It is shown by the solid line in the figure.
We note that the distribution function is not centered on $`\beta =0`$. The displacement of the center by about 9 mrad is much more than the probable error of $`\beta `$. We believe that it is caused by the effect of the Coriolis force on the LSC that will be discussed in more detail elsewhere \[\[Brown et al. (2005c)\]\].
## 7 Tilt-angle dependence of the center temperature
We saw from Fig. 11 that the increase of $`R_e`$ with $`\beta `$ led to an increase of the amplitude $`\delta `$ of the azimuthal temperature variation at the horizontal mid-plane. An additional question is whether the tilt-angle effect on this system has an asymmetry between the top and bottom that would lead to a change of the mean center temperature $`T_c`$ (see Eq. 4). \[Chillร et al.(2004)\] report such an effect for their $`\mathrm{\Gamma }=0.5`$ sample. For a Boussinesq sample with $`\beta =0`$ we expect that $`T_c=T_m`$ with $`T_m=(T_t+T_b)/2`$ ($`T_t`$ and $`T_b`$ are the top and bottom temperatures respectively), or equivalently that $`\mathrm{\Delta }_t=T_cT_t`$ is equal to $`\mathrm{\Delta }_b=T_bT_c`$. A difference between $`\mathrm{\Delta }_b`$ and $`\mathrm{\Delta }_t`$ will occur when the fluid properties have a significant temperature dependence \[\[Wu & Libchaber(1991), Zhang et al.(1997)\]\], i.e. when there are significant deviations from the Boussinesq approximation. For the sequence of measurements with the large apparatus and $`R=9.43\times 10^{10}`$ as a function of $`\beta `$ the mean value of $`\mathrm{\Delta }T=T_bT_t`$ was $`19.808\pm 0.018^{}`$C and $`T_cT_m`$ was 0.97C, indicating a significant non-Boussinesq effect. In Fig. 17a we show $`\mathrm{\Delta }_t`$ and $`\mathrm{\Delta }_b`$ as a function of $`\beta `$. One sees that increasing $`\beta `$ does not have a significant effect for our $`\mathrm{\Gamma }=1`$ sample. This is shown with greater resolution in Fig. 17b where $`T_cT_m=(\mathrm{\Delta }_t\mathrm{\Delta }_b)/2`$ is shown. We believe that the small variation, over a range of about $`5\times 10^3^{}`$C, is within possible systematic experimental errors and consistent with the absence of a tilt-angle effect.
## 8 Summary
In this paper we reported on an experimental investigation of the influence on turbulent convection of a small tilt angle $`\beta `$ relative to gravity of the axes of two cylindrical Rayleigh-Bรฉnard samples. The aspect ratios were $`\mathrm{\Gamma }1`$.
Where there was overlap, there were significant differences between our results and those obtained by \[Chillร et al.(2004)\] for a $`\mathrm{\Gamma }=0.5`$ sample. We found our system to establish a statistically stationary state quickly, within a couple of hours, after a Rayleigh-number change whereas \[Chillร et al.(2004)\] found long transients that they attributed to changes of the LSC structure. We found a very small depression of the Nusselt number $`๐ฉ`$ with increasing $`\beta `$, by about 4% per radian at small $`\beta `$. \[Chillร et al.(2004)\] found a decrease by 200% per radian for their sample.
In contrast to the very small effect of $`\beta `$ on $`๐ฉ`$, we found an increase of the Reynolds number $`R_e`$ by about 180% per radian for small $`\beta `$. The small effect on $`๐ฉ`$ in the presence of this large change of $`R_e`$ indicates that the heat transport does not depend strongly on the speed of the LSC sweeping over the boundary layers. Instead, $`๐ฉ`$ must be determined by instability mechanisms of the boundary layers, and the associated efficiency of the ejection of hot (cold) volumes (so-called โplumesโ) of fluid from the bottom (top) boundary layer.
It is interesting to note that the strong dependence of $`R_e`$ on $`\beta `$ in the presence of only a very weak dependence of $`๐ฉ`$ on $`\beta `$ can be accommodated quite well within the model of \[Grossmann & Lohse (2002)\]. The Reynolds number can be changed by introducing a $`\beta `$-dependence of the parameter $`a(\beta )`$ in their Eqs. (4) and (6). As pointed out by them, a change of $`a`$ has no influence on the predicted value for $`๐ฉ`$.
We also measured the frequency of rapid LSC reorientations that are known to occur for $`\beta =0`$. We found that such events are strongly suppressed by a finite $`\beta `$. Even a mild breaking or the rotational invariance, corresponding to $`\beta 0.04`$, suppresses re-orientations almost completely.
## 9 Acknowledgment
We are grateful to Siegfried Grossmann and Detlef Lohse for fruitful exchanges. This work was supported by the United States Department of Energy through Grant DE-FG02-03ER46080.
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# Spin-wave softening and Hundโs coupling in ferromagnetic manganites
## Abstract
Using one-orbital model of hole-doped manganites, we show with the help of Holstein-Primakov transformation that finite Hundโs coupling is responsible for the spin-wave softening in the ferromagnetic $`B`$-phase manganites. We obtain an analytical result for the spin-wave spectrum for $`J_\mathrm{H}t`$. In the limit of infinte Hundโs coupling, the spectrum is the conventional nearest-neighbor Heisenberg ferromagnetic spin-wave. The $`o(t/J_\mathrm{H})`$-order correction is negative and thus accounts for the softening near the zone boundary.
The observations of large magnetoresistence (LMR) in Nd<sub>0.5</sub>Pb<sub>0.5</sub>MnO<sub>3</sub> , giant magnetoresistence (GMR) and colossal magnetoresistence(CMR) in manganites ($`\mathrm{R}_{1x}\mathrm{A}_x\mathrm{MnO}_3`$, R is a rare earth element and A a divalent alkaline-earth metal) a decade ago s1 have rekindled much interest in these materials which have been known for half a centurys2 . Upon doping, the manganites undergo complicated transitions resulting in various magnetic, charge-ordering and orbital-ordering phases, showing the interplay between relevant spin, charge and orbital degrees of freedom. In particular, magnetism and electronic transport are clearly correlated. So it is widely believed that knowledge of spin dynamics can provide important information of the underlying physics of CMR. Perring et al first measured the spin waves in $`\mathrm{La}_{0.7}\mathrm{Pb}_{0.3}\mathrm{MnO}_3`$ for a broad range of $`q`$s3 . The magnon spectrum is well defined at low temperatures and can be accounted for by the nearest neighbor Heisenberg model. Subsequent measurements for $`\mathrm{Pr}_{0.63}\mathrm{Sr}_{0.37}\mathrm{MnO}_3`$ and $`\mathrm{Nd}_{0.7}\mathrm{Sr}_{0.3}\mathrm{MnO}_3`$s4 , $`\mathrm{Nd}_{0.7}\mathrm{Ba}_{0.3}\mathrm{MnO}_3`$ s5 further showed that the magnon spectrum deviates from the Heisenberg model and becomes softened near the zone boundary. So the behavior seems a universal phenomenon of manganites.
As is well known, a number of interactions such as spin-orbital coupling, Hundโs coupling, antiferromagnetic coupling between core spins, Coulomb interaction and dynamic Jahn-Teller effect coexist in manganites. These interactions are supposed to explain the existence of different phases of doped manganites. To explain the spin wave softening, various mechanisms were proposed. The authors of s4 further showed that the experimental spectrum can be reproduced reasonably well by an extended Heisenberg model. Furukawa s6 argued that the softening seems to be explainable by ferromagnetic Kondo lattice model with bandwidth narrower than the Hundโs coupling. Solovyev et als7 showed that the spin-wave behavior near the zone boundary has a purely magnetic spin origin, and neither the lattice deformation nor the orbital ordering are required to account for the softening. Dai et al argued that the observed magnon softening and broadening are due to strong magnetoelastic interactionss8 . And this magnon-phonon coupling was later treated quantitatively ins9 . Using ferromagnetic Kondo lattice model and composite operator method, Mancini et al obtained the softening spectrums10 . Shannon et al constructed a theory of spin wave excitations in the bilayer manganite $`\mathrm{La}_{1.2}\mathrm{Sr}_{1.8}\mathrm{Mn}_2\mathrm{O}_7`$ based on the simplest double exchange model and explained partly the softening behaviors11 . Krivenko et al showed that the scattering of spin excitations by low-lying orbital modes may cause the magnon softenings12 .
In this paper, we show that in the hole-doped manganites, the softening behavior might be of a purely electronic origin ,i.e., a strong but finite Hundโs coupling between the $`e_g`$ electron and the core spin. Since in the hole-doped manganite there is less than one $`e_g`$ electron per site on average and the $`d_{x^2y^2}`$ orbital energy is significantly higher than that of $`d_{3z^2r^2}`$s13 due to Jahn-Teller splitting, one orbital description is a reasonable approximation. As in s14 , we adopt the model Hamiltonian
$$H=t\underset{๐ข,๐ฃ}{}\underset{\sigma }{}c_{๐ข\sigma }^{}c_{๐ฃ\sigma }J_\mathrm{H}\underset{๐ข}{}๐ฌ_๐ข๐_๐ข+J_{\mathrm{AF}}\underset{<๐ข,๐ฃ>}{}๐_๐ข๐_๐ฃ\mu \underset{๐ข\sigma }{}c_{๐ข\sigma }^{}c_{๐ข\sigma }+U\underset{๐ข}{}n_๐ขn_๐ข$$
(1)
where $`t`$ is the double exchange hopping, $`๐ข,๐ฃ`$ are nearest sites, $`\mu `$ is the chemical potential for the fermions, $`c_{๐ข\sigma }`$ represents the $`e_g`$ electrons, $`J_\mathrm{H}`$ is the Hundโs coupling between the $`e_g`$ spin $`๐ฌ_๐ข=\frac{1}{2}c_๐ข^{}๐c_๐ข`$ and the the core spin $`๐_๐ข`$. $`J_{\mathrm{AF}}`$ is the antiferromagnetic interaction between the core spins, which is necessary to account for the $`G`$-phase parent $`(x=1)`$ manganites. The last term is the Hubbard Coulomb interaction. We use the Holstein-Primakoff transformation for the core spins ($`S=3/2)`$; $`S_๐ข^+=(2Sa_๐ข^{}a_๐ข)^{1/2}a_๐ข,S_๐ข^{}=a_๐ข^{}(2Sa_๐ข^{}a_๐ข)^{1/2},S_๐ข^z=Sa_๐ข^{}a_๐ข`$ and take the approximation $`(2Sa_๐ข^{}a_๐ข)^{1/2}(2Sa_๐ข^{}a_๐ข)^{1/2}`$. Homogeneity implies that $`a_๐ข^{}a_๐ข=a^{}a`$. Because we consider the low temperature case, we can drop the magnon quadratic term, hence the total Hamiltonian can be written as
$`H`$ $`=`$ $`t{\displaystyle \underset{๐ข,๐ฃ}{}}{\displaystyle \underset{\sigma }{}}c_{๐ข\sigma }^{}c_{๐ฃ\sigma }\mu {\displaystyle \underset{๐ข\sigma }{}}c_{๐ข\sigma }^{}c_{๐ข\sigma }+U{\displaystyle \underset{๐ข}{}}n_๐ขn_๐ข{\displaystyle \frac{1}{2}}J_\mathrm{H}A{\displaystyle \underset{๐ข}{}}(s_๐ข^+a_๐ข^{}+s_๐ข^{}a_๐ข)J_\mathrm{H}S{\displaystyle \underset{๐ข}{}}s_๐ข^z`$ (2)
$`+A^2J_{\mathrm{AF}}{\displaystyle \underset{๐ข,๐ฃ}{}}a_๐ขa_๐ฃ^{}+J_{\mathrm{AF}}ZNS^22ZJ_{\mathrm{AF}}S{\displaystyle \underset{๐ข}{}}a_๐ข^{}a_๐ข+J_\mathrm{H}{\displaystyle \underset{๐ข}{}}s_๐ข^za_๐ข^{}a_๐ข`$
, where $`A^2=2Sa^{}a,Z=6`$ is the coordination number of the core spins. To use the composite operator method, we consider the doublet $`B(๐ข)=(\begin{array}{cc}a_๐ข,& s_๐ข^+\end{array})^T`$. The equation of motion for $`B(๐ข)`$ is
$$i_tB(๐ข)=[B(๐ข),H]=\left(\begin{array}{c}\frac{1}{2}J_\mathrm{H}As_๐ข^++A^2J_{\mathrm{AF}}_๐a_{๐ข+๐}2ZJ_{\mathrm{AF}}Sa_๐ข+J_\mathrm{H}s_๐ข^za_๐ข\\ t_๐(c_๐ข^{}c_{๐ข+๐}c_{๐ข+๐}^{}c_๐ข)J_\mathrm{H}As_๐ข^za_๐ข+J_\mathrm{H}Ss_๐ข^+J_\mathrm{H}s_๐ข^+a_๐ข^{}a_๐ข\end{array}\right)$$
(3)
Composite operator method assumes that the right-hand side can be expressed as
$$[B(๐ข),H]=\underset{๐ฃ}{}\epsilon (๐ข,๐ฃ)B(๐ฃ)$$
(4)
with $`\epsilon (๐ข,๐ฃ)`$ determined in the following way,
$$\epsilon (๐ข,๐ฃ)=\underset{๐ฅ}{}m(๐ข,๐ฅ)I^1(๐ฅ,๐ฃ)$$
(5)
where $`I(๐ข,๐ฃ)=[B(๐ข),B^{}(๐ฃ)],m(๐ข,๐ฃ)=[i_tB(๐ข),B^{}(๐ฃ)]`$, and $``$ represents the expectation value. Thus $`\epsilon (๐ข,๐ฃ)`$ contains some parameters to be determined self-consistently. This approach was proposed for Hubbard model originallys14 and recent intensive studiess15 show credible agreement with Monte Carlo method. In our case (again due to homogeneity, $`s_๐ข^z=s^z`$)
$`I(๐ข,๐ฃ)`$ $`=`$ $`\delta _{\mathrm{๐ข๐ฃ}}\mathrm{diag}(1,2s^z)`$
$`m_{11}(๐ข,๐ฃ)`$ $`=`$ $`\delta _{\mathrm{๐ข๐ฃ}}(J_\mathrm{H}s^z2ZSJ_{\mathrm{AF}})+A^2J_{\mathrm{AF}}{\displaystyle \underset{๐}{}}\delta _{๐ฃ,๐ข+๐}`$
$`m_{12}(๐ข,๐ฃ)`$ $`=`$ $`\delta _{\mathrm{๐ข๐ฃ}}J_\mathrm{H}(As^zs_๐ข^{}a_๐ข)`$
$`m_{22}(๐ข,๐ฃ)`$ $`=`$ $`tp_1{\displaystyle \underset{๐}{}}(\delta _{\mathrm{๐ข๐ฃ}}\delta _{๐ฃ,๐ข+๐})+J_\mathrm{H}As_๐ข^{}a_๐ข`$
$`+2J_\mathrm{H}Ss^z2J_\mathrm{H}s_๐ข^za_๐ข^{}a_๐ข`$
where $`p_1=_\sigma c_{๐ข\sigma }^{}c_{๐ข\sigma },p_2=s_๐ข^{}a_๐ข,p_3=s_๐ข^za_๐ข^{}a_๐ข`$. In the $`๐ค`$-space
$`m_{11}(๐ค)`$ $`=`$ $`(J_\mathrm{H}s^z2ZSJ_{\mathrm{AF}})+ZA^2J_{\mathrm{AF}}\gamma _๐ค`$
$`m_{12}(๐ค)`$ $`=`$ $`J_\mathrm{H}(As^zp_2)`$
$`m_{22}(๐ค)`$ $`=`$ $`J_\mathrm{H}Ap_2tZp_1(1\gamma _๐ค)+2J_\mathrm{H}Ss^z2J_\mathrm{H}p_3`$
We assume that at $`T=0`$K,$`a^{}a=0`$, which satisfies self-consistency using the resulting retarded Greenโs function and spectral theorem. Then condition $`\omega _{|๐ค=0}=0`$ requires that $`p_3=\frac{1}{2}p_2(A+\frac{p_2}{s^z})`$. So the $`\epsilon `$-matrix is
$`\epsilon _{11}(๐ค)`$ $`=`$ $`J_\mathrm{H}s^z2ZSJ_{\mathrm{AF}}(1\gamma _๐ค)`$
$`\epsilon _{12}(๐ค)`$ $`=`$ $`{\displaystyle \frac{J_\mathrm{H}}{2s^z}}(As^z+p_2)`$
$`\epsilon _{21}(๐ค)`$ $`=`$ $`J_\mathrm{H}(As^z+p_2)`$
$`\epsilon _{22}(๐ค)`$ $`=`$ $`J_\mathrm{H}S{\displaystyle \frac{tZp_1}{2s^z}}(1\gamma _๐ค)+J_\mathrm{H}A{\displaystyle \frac{p_2}{s^z}}+J_\mathrm{H}{\displaystyle \frac{p_2^2}{2s^z^2}}`$
and the Greenโs function is
$`D_{11}(\omega ,๐ค)`$ $`=`$ $`{\displaystyle \frac{\omega (J_\mathrm{H}S\frac{tZp_1}{2s^z}(1\gamma _๐ค)+J_\mathrm{H}A\frac{p_2}{s^z}+J_\mathrm{H}\frac{p_2^2}{2s^z^2})}{(\omega \omega _1(๐ค))(\omega \omega _2(๐ค))}}`$
$`D_{12}(\omega ,๐ค)`$ $`=`$ $`D_{21}(๐ค)={\displaystyle \frac{J_\mathrm{H}(As^z+p_2)}{(\omega \omega _1(๐ค))(\omega \omega _2(๐ค))}}`$
$`D_{22}(\omega ,๐ค)`$ $`=`$ $`{\displaystyle \frac{2s^z[\omega (J_\mathrm{H}s^z2ZSJ_{\mathrm{AF}}(1\gamma _๐ค))}{(\omega \omega _1(๐ค))(\omega \omega _2(๐ค))}}`$
where $`\omega _{1,2}(๐ค)`$ are acoustical and optical branches of the spin excitations. Using
$$p_2=\frac{1}{N}\underset{๐ค}{}\frac{i}{2\pi }๐\omega \underset{\eta 0}{lim}\frac{D_{12}(\omega +i\eta ,๐ค)D_{12}(\omega i\eta ,๐ค)}{e^{\beta \omega }1}$$
(6)
we have at $`T=0`$K, $`p_2=0`$, therefore, $`p_3=0`$. Accordingly, in this scheme,there are two parameters left: $`s^z,p_1`$. The acoustical magnon spectrum can be expanded as a Taylor series which manifests the role of Hundโs coupling
$$\omega _1(๐ค)=t\underset{n=0}{\overset{\mathrm{}}{}}(\frac{t}{J_\mathrm{H}})^na_n(1\gamma _๐ค)^{n+1}$$
(7)
where as usual,$`\gamma _๐ค=Z^1_๐e^{i๐ค๐}`$. The first few $`a_n`$ are
$`a_0`$ $`=`$ $`{\displaystyle \frac{3(4S^2\frac{J_{\mathrm{AF}}}{t}+p_1)}{s^z+S}}`$
$`a_1`$ $`=`$ $`{\displaystyle \frac{9S(4S\frac{J_{\mathrm{AF}}}{t}s^zp_1)^2}{s^z(s^z+S)^3}}`$
$`a_2`$ $`=`$ $`{\displaystyle \frac{27(4S\frac{J_{\mathrm{AF}}}{t}s^zp_1)^3S(Ss^z)}{s^z^2(S+s^z)^5}}`$
$`a_3`$ $`=`$ $`{\displaystyle \frac{81(S^23Ss^z+s^z^2)(4S\frac{J_{\mathrm{AF}}}{t}s^zp_1)^4S}{s^z^3(S+s^z)^7}}`$
In the small $`k`$ limit, $`\omega Dk^2,D=(4S^2\frac{J_{\mathrm{AF}}}{t}+p_1)/(s^z+S)`$. Note that the hopping energy $`tp_1`$ is negative and when it overcomes the AF term, the resulting magnon stiffness $`D`$ is positive. Our numerical results show that this self-consistency is satisfied. Expression (7) suggests that the softening comes from the finite $`J_\mathrm{H}`$. To fix the parameters $`s^z`$, we use spectral theorem and get $`s^z=\frac{1}{2}(1x)`$, where $`x`$ is the dopant concentration. To fix $`p_1`$, we need the fermion sector. Using the notations in s10 for the fermion operator $`\psi (๐ข)=(\begin{array}{c}\xi _๐ข,\eta _๐ข,\xi _๐ข,\eta _๐ข\end{array})^T`$, where $`\xi _\sigma =(1n_\sigma )c_\sigma ,\eta _\sigma =n_\sigma c_\sigma `$ are the Hubbard operators, we obtain the retarded Greenโs function for $`\psi `$ in the large$`U`$ limit at zero temperature.
$$G^R(\omega ,๐ค)=\mathrm{diag}(\frac{1}{\omega E_1(๐ค)},0,\frac{x}{\omega E_3(๐ค)},\frac{1x}{\omega E_4(๐ค)})$$
(8)
(โdiagโ means diagonal matrix) with $`E_1(๐ค)=\mu +6t\gamma _๐ค\frac{1}{2}SJ_\mathrm{H},E_2(๐ค)=U\mu +6tu+6tv\gamma _๐ค,E_3(๐ค)=[24tx\gamma _๐ค2\mu x+SJ_\mathrm{H}x+12tp_{}\gamma _๐ค12t\gamma _๐ค+12t\mathrm{\Delta }_{}]/(2x),E_4(๐ค)=[2Ux+2\mu xSJ_\mathrm{H}x+2U+12tp_{}\gamma _๐ค+SJ_\mathrm{H}2\mu +12t\mathrm{\Delta }_{}]/[2(1x)]`$ , where $`\mathrm{\Delta }`$ is related to the nearest-neighbor correlations of the Hubbard operators : $`\mathrm{\Delta }_\sigma =\xi _\sigma (๐ข+๐)\xi _\sigma ^{}(๐ข)\eta _\sigma (๐ข+๐)\eta _\sigma ^{}(๐ข)`$. In this scenario, $`E_1,E_3`$ are partially filled and $`E_2,E_4`$ are empty. The relevant parameters are $`\mu ,\mathrm{\Delta }_{},p_{}`$. We have three equations to fix them $`1x=2C_{11}^FC_{22}^FC_{33}^FC_{44}^F,,\mathrm{\Delta }_{}=C_{11}^{F\gamma },C_{11}^F=C_{33}^F`$ , where $`C^F=\psi (๐ข)\psi ^{}(๐ข),C^{F\gamma }=\psi (๐ข+๐)\psi ^{}(๐ข)`$. We know that $`C_{22}^F=0`$ and $`C_{44}^F=1x`$. Thus $`C_{11}^F=x=C_{33}^F`$, so $`E_3`$ is empty, i.e., only $`E_1`$ is partially filled. Hence only $`\mu `$ is relevant to our problem and can be fixed by $`x=N^1_๐ค\theta (E_1(๐ค))`$ where $`\theta (x)`$ is the usual step function. The hopping energy is
$$tp_1=tC_{11}^{F\gamma }=\frac{t}{N}\underset{๐ค}{}\theta (E_1(๐ค))\gamma _๐ค<0$$
(9)
The other two parameters $`\mathrm{\Delta }_{},p_{}`$ can also be determined by $`t\mathrm{\Delta }_{}=tp_1>0,24tx+12t(p_{}1)=2\mu x+SJ_\mathrm{H}x+12t\mathrm{\Delta }_{}`$. Further analysis show that for $`J_\mathrm{H}>2.5t`$, the whole scheme is self-consistent. Fig.1 shows that two relevant fermion bands for $`x=0.301,t=1,J_\mathrm{H}=3.0`$ ( in unit of $`t`$).
It is seen from the magnon spectrum (7) that we can estimate the two model parameters $`t`$ and $`J_\mathrm{H}`$ from measured data. Fig.2 shows the comparison between our calculated result for the prescribed antiferromagnetic coupling $`J_{\mathrm{AF}}=0.01`$ and the measured result at $`T=10`$K for $`\mathrm{Pr}_{0.63}\mathrm{Sr}_{0.37}\mathrm{MnO}_3`$ ins4 . The solid curve in the left panel is the fit to a nearest -neighbor Heisenberg model and gives the value at zone boundary about 34.2meV. This corresponds to the the uppermost curve in the right panel. The comparison gives the hopping energy $`t0.462`$eV. The circles are the data measured and give the value at zone boundary about 23meV, corresponding to the point 0.05 in the right panel. This point corresponds to $`J_\mathrm{H}=3.2t1.48`$ eV. Note that the ratio $`J_\mathrm{H}/t`$ is very close to the values of interaction from a number of referencess16 s17 . It is worth noting that the nearest-neighbor Heisenberg interaction alone can not account for the Curie temperature. The fitting curve in the left panel corresponds to the nearest-neighbor Heisenberg spectrum $`\omega (๐ค)51.3(1\gamma _๐ค)`$ meV. In the mean field theory,the Curie temperature $`T_c`$ corresponding to the spectrum $`\omega _๐ค=2ZS^{}J(1\gamma _๐ค)`$ is $`k_BT_c=\frac{2}{3}JZS^{}(S^{}+1)`$ (here $`S^{}=S+\frac{1}{2}(1x)`$ is the effective spin). This gives $`T_C^{\mathrm{MF}}500`$K. Taking into account that in three dimensions for a simple cubic lattice , the real Curie temperature $`T_C`$ and $`T_C^{\mathrm{MF}}`$ have relations18 : $`T_C=0.75T_C^{\mathrm{MF}}`$, we get $`T_c375`$ K, higher than the real value 315 K.
To conclude this paper, we present some discussions and comments. In the derivation of the series expression of the magnon spectrum, we have used the approximation $`(2Sa_๐ข^{}a_๐ข)^{1/2}(2Sa_๐ข^{}a_๐ข)^{1/2}`$ in the Holstein-Primarkov transformation. This can be satisfied at very low temperatures. Further, the quartic term $`J_{\mathrm{AF}}_{๐ข,๐ฃ}a_๐ข^{}a_๐ขa_๐ฃ^{}a_๐ฃ`$ is neglected because $`J_{\mathrm{AF}}`$ is very small and the magnon fluctuation at zero temperature is negligible. The series expression (7) of the accoustic magon dispersion shows alternating behavior; convergence is guaranteed when $`J_\mathrm{H}S>3`$. The model parameters $`t,J_\mathrm{H},J_{\mathrm{AF}}`$ and the hopping energy $`p_1`$ can be estimated by fitting experimental data. There is a simple physical picture for the deviation of the magnon spectrum from that of Heisenberg model. The interaction between core spins is induced by the hopping of $`e_g`$-electrons and the dominant term is linear in $`t`$. If the Hundโs coupling $`J_\mathrm{H}`$ is infinite,only the dominant term plays the role. The $`e_g`$ electron and core spin must add up to a total spin-2 to minimize the energy in the $`B`$-phase. So the actual background for spin excitation is just that in the simple Heisenberg model. But for finite $`J_\mathrm{H}`$, high orders of the mediated interaction between core spins make some difference. Our result (7) agrees with the conclusions from random phase approximations19 , which provides an integral equation for the dispersion relation. The strength of induced ferromagnetic interaction is determined by the hopping energy of the conduction fermions. For the approach presented to be self-consistent, the spin-wave stiffness must be positive. The ferromagnetic order becomes unstable at certain filling when the the stiffness vanishes. Though zone boundary spin wave softening can be explained by the spin dynamics in the Kondo lattice model, as shown in this paper. The origin of the behavior is still an issue of debate. Based on the observed proximity of phonon dispersion and magnon dispersion and the anisotropic spin-wave broadening, Dai et als8 argue that strong magnon-phonon coupling is needed for a complete understanding of the low temperature spin dynamics of manganites. Quite recently, Endoh et al concludeds20 that the ferromagnetic magnons in Sm<sub>0.55</sub>Sr<sub>0.45</sub>MnO<sub>3</sub> is of orbital nature since the magnon dispersion shows anisotropy which is mainly determined by the short range correlation of the $`e_g`$ orbitals. They explained that the anisotropic magnon dispersion are attributed to long range magnetic interactions based on fitting the data to Heisenberg model with long range interactions. We believe that if orbital degrees of freedom are taken into account in our model, the resulting magnon spectrum will be anisotropic since orbital degrees of freedom bring anisotropy into the system. Finally, we remark that as manganites are very complex systems, there might be multiple mechanisms contributing to a single phenomenon. The analysis provided in this paper shows that Hund coupling might be of primary importance.
SSF is grateful to Prof. P. Schlottmann, Prof. Mancini, and Dr. Avella for their helpful discussions. This work is supported in part by the Army High Performance Computing Research Center(AHPCRC) under the auspices of the Department of the Army, Army Research Laboratory (ARL) under Cooperative Agreement number DAAD19-01-2-0014.
Corresponding author,e-mail:shixiang.feng@famu.edu
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# Interaction between Physics and Cosmology
## Abstract
Recent results indicate the presence of a cosmological constant (or related dark energy) in the universe. It has been conjectured recently that the interaction parameters of physical theories may be dependant on the size parameter of the universe, related to the cosmological constant. We investigate whether such effects will help in explaining baryogenesis in early universe. They do seem to succeed.
## I. Introduction
Our knowledge of cosmology has improved remarkably in recent years. along with it has come many new or revived ideas and concepts. Recent observation of accelerated expansion of the universe requires the presence of a cosmological constant, introduced and withdrawn by Einstein long ago. The cosmological constant is related to the enegy of vacuum. Its small value is difficult to understand in particle physics.The cosmological constant leads to a de Sitter universe with a finite size; the larger the constant,the smaller the universe. Bjorken has extended this idea of a relation between size of universe and the vacuum energy to other physical parameters which were earlier taken to be constant. During the period of inflationary expansion of the universe the vacuum energy is different and so is the size of the universe. We expect physical interaction parametrs to be also different.Inflationary expansions may occur during phase transitions in early universe.Baryogenesis in early universe is expected to occur during such phase transitions.If physical parameters change, as suggested by Bjorken, then this will affect baryogenesis. We shall consider this effect in this paper.
## II. Variaton of Standard Model Parameters
The size of the universe R and the value of cosmological constant $`\mathrm{\Lambda }`$ are related by the equation $`\mathrm{\Lambda }/3=1/R_{\mathrm{}}^2=8\pi G\rho /3=H_{\mathrm{}}^2`$ where the subscript refers to the values in the limit of very large times. Bjorken assumed that โAll dimensionful parameters X of the standard model may vary with $`R_{\mathrm{}}`$ but that to leading approximation they are straight lines in a log- log plot, i.e. they satisfy a simple renormalisation group equationโ. The equation satisfied is of the form:
$`R^2\frac{X}{R^2}=1/2\mu \frac{X}{\mu }=p_XX+\mathrm{}..`$ .
where we have dropped the subscript on R.
We consider the behaviour of $`\mathrm{\Lambda }_{QCD}`$, the cut of mass in QCD, v,the Vacuum expectation value of the Higgs field in electroweak theory and $`M_{Gut}`$, the unification energy in grand unified theories. The parameters $`\mathrm{\Lambda }_{QCD},v,M_{Gut}`$ all vary with R. This behaviour is shown in Fig.1. According to present ideas these parameters are all constant and never go up to $`M_{Pl}`$ as shown in the figure.
As $`1/(\alpha _s(q^2))=b_sln(q^2/(\lambda ^2))`$
and $`\mathrm{\Lambda }_{QCD}^2=(M_{Pl}^2R^2)^{p_s}M_{Pl}^2`$
with $`p_s1/3`$, from Fig.1, it follows that
$`1/(\alpha _s(q^2,R^2))=b_sp_sln(M_{Pl}^2R^2)b_sln(M_{Pl}^2/(q^2))`$ These equations imply a renormalisation group equation for the coupling constant $`\alpha _s`$ of the form :
$`R^2\frac{}{R^2}(1/(\alpha _s))=b_sp_s+O(\alpha _s)`$.
This implies a logarithmic increase for the coupling constant as R decreases. When $`R\mathrm{}`$ the coupling constant vanishes. In an infinite universe the standard model trivialises to a free field theory. All interacrions depend on the existence of a boundary to the universe. In Planck limit fields becomes strongly coupled. The masses of quarks are given by $`m=gv`$. So $`m/v`$ increases logrithmically as R decreases.
## III. Baryon Assymmetry in the Universe
There are many scenarios for baryogenesis. We shall confine ourselves to the one that occurs at the time of electroweak phase transition . This scenario uses only parameters of the standard model known from experiments.In this model the baryon aymmetry of the universe (BAU) is given by
$`BAU=J[(m_t)^2(m_u)^2)((m_t)^2(m_c)^2)((m_c)^2(m_u)^2)((m_b)^2(m_s)^2)((m_b)^2(m_d)^2)((m_s)^2(m_d)^2)/(T^{12}]`$.
$`J=\mathrm{sin}(\theta _{12})\mathrm{sin}(\theta _{13})\mathrm{sin}(\theta _{23})\mathrm{sin}(\delta _{CP})`$.
Here J is the CP violating factor determined from decays of K and now B mesons. The angles are mixing angles of CKM matrix for quark mixing .
This BAU has a value of $`10^{21}`$ while the required value based on nucleosynthesis in early universe is of the order of $`10^{10}`$
The new point we are trying to make is that the value of R at the time of phase transition is different due to supercooling and inflationary expansion. The value of R is 1 cm at the e.w. phase transition, while the value of R for our universe is $`10^{28}cm`$. The physical parameters of the standard model change with R as discussed in the previous section. This will give an additional factor of $`(log(R/R_{ew})^{12}`$ in BAU. The value of this factor is $`10^{16}`$. So the BAU is about $`10^5`$ which is a much more reasonable value than the earlier one, as detailed analysis will bring it down by a few orders.
There is thus hope that baryogenesis scenario at the time of elactroweak phase transition may explain the value of BAU that is required .
## IV Discussion
As mentioned earlier there are many scenarios for baryogenesis .We have considered the one closest to known physics in the laboratory. This scenario has two major difficulties. The first is the requirement of a small mass for the Higgs particle to ensure first order phase transition. This seems already inconsistent with the observations. To study this difficulty in the present framework, the constraints on the mass of the Higgs have to be reworked taking the variation of the parameters at the time of phase transition. This requires detailed calculations along the lines done in our earlier work. The second difficulty is the very small value of BAU that is predicted. We have shown above that the second difficulty seems to be overcome if parameters depend on the size of the universe.
The simplest extension of our scenario involves minimal super- symmetry with an extra Higgs particle. This has to await signals of super- symmetry in the laboratory. The other scenarios assume fields whose parameters ( even their very existence) is unobserved and unconfirmed in thelaboratory. The scenarios based on string theories and brane models are again in a different realm.
## V References
1. J.D.Bjorken, Phys. Rev. D 67,043508-18 (2003).
2. A.S.Majumdar,S.K.Sethi,S.Mahajan, A.Mukherjee, N.Panchapakesan,and R.P.Saxena, Mod. Phys. Lett., A9,459 (1994). G.R.Farrar,and M.E.Shaposhnikov Phys.Rev. Lett.,70,2833 (1993), 71,210(E) (1993), Phys.Rev. D50, 774 (1994).
3.M.Trodden, Rev. Mod. Phys. (1999), A.Riotto and M.Trodden,hep-phys/9901362, M.Quiros, hep-phys/9901312.
4. For a similar approach but with a different motivation, see M.Berkooz, Y.Nir, and T. Volansky, Phys. Rev. Lett. 93,051301-4 (2004).
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# Bose-Einstein Condensation Picture of Superconductivity in "YBa"โโข"Cu"โโข"O"โโข(91โข"K") and "YBa"โโข"Cu"โโข"Se"โโข(371โข"K"). (Dilute metals).
## Abstract
A metal dilution degree in the compounds $`\text{YBa}_2\text{Cu}_3\text{O}_7`$ and $`\text{YBa}_2\text{Cu}_3\text{Se}_7`$ is defined as $`z=(r_{\text{Me}}/r_{\text{O}2;\text{Se}2})^3.`$ A substitution of oxygen by selenium changes z by $`8`$ times and the Bose-Einstein condensation temperature equals $`T_{\text{cSe}}=T_{\text{cO}}8^{2/3}=364\text{K}.`$ The โactiveโ electron pairs density in $`\text{YBa}_2\text{Cu}_3\text{O}_7`$ is about $`1.7\times 10^{20}\text{cm}^3.`$ The electron effective mass is about $`5m_e`$ and is proportional to the dielectric constant.
The superconductivity at $`371\text{K}`$ in $`\text{YBa}_2\text{Cu}_3\text{O}_7`$ was found measuring the magnetic properties in bib1 . However, the substitution of the oxygen ions $`(r_{02}0.5\AA )`$ by the selenium ones $`(r_{\text{Se}2}1.0\AA )`$ made the lattice unstable and the compound irreproducible.
Densities of the electron pairs, transition temperatures and effective masses at the โphysicalโ dilution of metals
in $`\text{Na}_{0.04}\text{NH}_3(T_c200\text{K},m^{}5m_e)`$ bib2 ,
in $`\text{Na}_{0.05}\text{WO}_3(T_c91\text{K},m^{}10m_e)`$ bib3 , and
in $`\text{Ag}_2(\text{Ag}_3\text{Pb}_2\text{H}_2\text{O}_6)(T_c400\text{K},m^{}7.5m_e)`$ bib4 are known and comply with the BEC model bib5 .
The โactiveโ electron pair density in oxides of the $`(1237)`$ type corresponds to the metal orbital occupancy, which depends upon the acceptor properties of the oxygen or selenium ions (the โchemicalโ dilution) bib6 . (The โactiveโ electron density in the Noble Gases condensates determines by interaction of the atom ground state orbitals bib6 ).
Let us assume that the Me and $`\text{O}^2`$ orbital occupancy are inversely proportional to the orbital volumes: $`z=(r_{\text{Me}}/r_{O2})^3;z=8.0`$ for the pair $`\text{Se}^2\text{O}^2.`$ Therefore, the pair density in YBaCuSe is 8 times larger, than in YBaCuO. For the BEC mechanism, $`T_{\text{cSe}}=T_{\text{cO}}8^{2/3}=914=364\text{K}`$ that complies with the data bib1 .
About one metal atom falls at one oxygen atom in YBaCuO. The pair number in the unit cell equals โ$`2.5.`$ Every electron pair occupies the volume $`69.4\AA ^3.`$ For the occupancy $`z_{\text{BaโO}}=(2.19/0.5)^3=84`$ we have for the โactiveโ electron pair density is about $`1.7\times 10^{20}\text{cm}^3.`$ The transition temperature $`T_{c\text{O}}=92\text{K}`$ for the BEC mechanism and the effective mass $`m^{}5m_e`$ (the dielectric constant $`20`$โโ).
The effective electron mass magnitudes are very close in all of these four systems bib1 ; bib2 ; bib3 ; bib4 and comply with the dielectric constants.
The compounds studied in bib1 and bib4 show the superconductivity of the โgossamerโ kind.
A use of hydrogen for blocking of a part of the oxygen valence can preserve the lattice structure at the electron concentration increase leading to an increase of $`T_{c\text{O}}`$ bib7 .
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# References
$``$ Introduction: Everyone who has taught introductory physics is familiar with the confusion students sometimes experience between velocity and acceleration. When a stone is thrown upwards in a uniform gravitational field its acceleration is a negative constant whereas its velocity decreases linearly from positive to negative. If the same process were viewed from the frame of an inertial observer initially moving upwards faster than the stone, the stoneโs velocity would be always negative but its acceleration would be unchanged. We shall argue that a similar confusion exists in identifying an invariant measure of the quantum gravitational back-reaction on inflation. What one wants to do is quantify the tendency for quantum processes to decelerate the universe, and this is independent of how rapidly the universe may be expanding with respect to some velocity field.
Simple arguments indicate that gravitational back-reaction should slow inflation, either in pure gravity or in scalar-driven inflation . The idea is that inflation rips from the vacuum a continuous stream of virtual long wavelength quanta which are massless and not conformally invariant. The gravitational interactions between these particles induces a negative energy density that gradually increases as more and more of these particles come into causal contact with one another.
Support for these ideas has come from detailed perturbative computations, at 2-loop order in $`\mathrm{\Lambda }`$-driven inflation and at 1-loop order in scalar-driven inflation . Unruh quite properly criticized these results on the grounds that they were obtained by regarding the expectation value of the gauge-fixed metric as an observable . One way of resolving his concern is to do the computation in other gauges. That was too difficult for the 2-loop gravitational process but a simple check in a completely different gauge showed no change in the claimed back-reaction for scalar-driven inflation .
A much better way of addressing Unruhโs criticism is to construct an invariant operator that measures the expansion rate even when perturbations are present. One can then compute the expectation value of this operator. A crude scalar measure of the spacetime expansion rate was obtained using the inverse conformal dโAlembertian acting upon unity . This operator can be promoted to an invariant be evaluating it at an physically defined observation point. When the 1-loop expectation value of that invariant was evaluated for a model which had previously seemed to show back-reaction at one loop, the result was zero secular back-reaction . Curiously, the scalar observable contributed nothing to this. When its expectation value was computed in the old, gauge-dependent coordinate systems it continued to show precisely the old back-reaction effect. The nullification of the effect was entirely due to the corrections needed to define the observation point invariantly.
A major step forward was taken by Geshnizjani and Brandenberger , who worked out how to apply the standard measure of expansion to the case of scalar-driven inflation. In this situation the scalar inflaton $`\phi (t,\stackrel{}{x})`$ provides a preferred 4-velocity field: <sup>1</sup><sup>1</sup>1Hellenic indices take on spacetime values while Roman indices take on space values. Our conventions are that the metric $`g_{\mu \nu }`$ has spacelike signature and the curvature tensor equals $`R_{\beta \mu \nu }^\alpha \mathrm{\Gamma }_{\nu \beta ,\mu }^\alpha +\mathrm{\Gamma }_{\mu \rho }^\alpha \mathrm{\Gamma }_{\nu \beta }^\rho (\mu \nu )`$.
$$u^\mu (t,\stackrel{}{x})\frac{g^{\mu \nu }(t,\stackrel{}{x})_\nu \phi (t,\stackrel{}{x})}{\sqrt{g^{\alpha \beta }(t,\stackrel{}{x})_\alpha \phi (t,\stackrel{}{x})_\beta \phi (t,\stackrel{}{x})}}.$$
(1)
In this expression we mean the quantum metric and the full inflaton operator, including its classical background part and the quantum perturbation:
$$\phi (t,\stackrel{}{x})=\phi _0(t)+\delta \phi (t,\stackrel{}{x}).$$
(2)
Expression (1) for $`u^\mu (t,\stackrel{}{x})`$ obviously gives a timelike unit vector which is well-defined for small perturbations. By taking one third of its divergence one obtains a scalar measure of the expansion rate. The scalar becomes a full invariant when evaluated at a physically defined observation point. Owing to the spatial homogeneity and isotropy of the state it is only necessary to invariantly fix the time. The preferred coordinate system is one in which the full scalar agrees with its background value. That is, one solves perturbatively for the operator $`\tau [\phi ,g](t,\stackrel{}{x})`$ which enforces the condition:
$$\phi (\tau (t,\stackrel{}{x}),\stackrel{}{x})=\phi _0(t).$$
(3)
The invariant expansion measure is:
$$(t,\stackrel{}{x})\frac{1}{3}u_{;\mu }^\mu (\tau (t,\stackrel{}{x}),\stackrel{}{x})=\frac{1}{3}\frac{1}{\sqrt{g}}_\mu \left(\sqrt{g}u^\mu \right)(\tau (t,\stackrel{}{x}),\stackrel{}{x}).$$
(4)
When the expectation value of (4) was evaluated for models which seemed to show secular back-reaction at 1-loop order, the result was no secular effect . Computations of quantities such as the stress tensor continue to show secular back-reaction at one loop when observables are evaluated at a gauge-fixed coordinate point instead of at the physically defined location (3). It has also been shown that secular back-reaction can seem to occur at 1-loop order in a two-scalar model if one uses the other, spectator, scalar to measure the expansion rate .
It seems clear that there is no secular one loop back-reaction in conventional scalar-driven inflation . <sup>2</sup><sup>2</sup>2But unconventional models might show it . This actually agrees with the putative physics of the effect. The causative agent is 1-loop particle production so it would have been fortuitous to see a secular gravitational response at the same order. Secular growth can only come as more and more of the produced particles come into interaction with one another and that must be delayed โ in perturbation theory โ until at least 2-loop order. Scalar models can be contrived which do show an invariantly quantified back-reaction effect at higher loop order . And the problem will not go away because of the enormous potential impact of back-reaction in a realistic model .
A particularly interesting model in which to investigate back-reaction is $`\mathrm{\Lambda }`$-driven inflation. The Lagrangian is just that of gravity with a positive cosmological constant:
$$=\frac{1}{16\pi G}\left(R2\mathrm{\Lambda }\right)\sqrt{g},$$
(5)
and the perturbative background would be de Sitter. One must certainly go to 2-loop order to see back-reaction in this model . On the other hand, the system is uncomplicated by the classical evolution of a scalar inflaton. The classical result for this model is no evolution at all so, if one sees evolution then it is due to quantum gravitational back-reaction; except that one still has to invariantly quantify the effect and there is no scalar inflaton to serve as a clock everyone can agree upon. Of course one could simply define an ersatz scalar any number of ways. For example, one could take the invariant volume of the past light-cone, measured from the observation point back to the initial value surface upon which the quantum state is released. This transforms as a scalar, and it certainly increases monotonically in the timelike direction. The trouble is that many such โscalarsโ can be defined and none of them has the privileged role that the scalar inflaton plays in setting the zero of time for the post-inflationary universe.
$``$ The de Sitter example: The ambiguity and the conflicting results it can produce are easy to understand on the classical level in de Sitter spacetime. Suppose we take the โscalarโ to be the time $`t_c`$ of the โclosedโ co-ordinate system: <sup>3</sup><sup>3</sup>3De Sitter spacetime has the topology of $`S^{D1}\times \mathrm{}`$ and it is natural to cover the full manifold by using a co-ordinate system in which the spatial sections are $`S^{D1}`$.
$$Closedds^2=dt_c^2+\frac{1}{H^2}\mathrm{cosh}^2(Ht_c)d\mathrm{\Omega }_3^2,$$
(6)
$$\mathrm{}<t_c<\mathrm{},0\alpha _i\pi (i=1,2),0\alpha _3<2\pi .$$
(7)
It is elementary to compute the velocity field (1) and the expansion rate (4) in this system:
$$Closedu_c^\mu =\delta _0^\mu H_c=H\mathrm{tanh}(Ht_c).$$
(8)
The integral curves of this velocity field are initially drawing together โ so the expansion rate starts negative โ but they eventually draw apart, resulting in positive expansion.
Now suppose we adopt as our โscalarโ the time $`t_o`$ of the โopenโ co-ordinates sub-manifold: <sup>4</sup><sup>4</sup>4The โopenโ co-ordinate system covers half of the full de Sitter manifold and has the topology of $`\mathrm{}^D`$.
$$Opends^2=dt_o^2+e^{2Ht_o}d\stackrel{}{x}^2.$$
(9)
$$\mathrm{}<t_o<\mathrm{},0\stackrel{}{x}<\mathrm{}$$
(10)
The resulting velocity field (1) and expansion rate (4) are:
$$Openu_o^\mu =\delta _0^\mu H_o=H.$$
(11)
For the โopenโ velocity field the expansion rate is a positive constant throughout the sub-manifold. The latter includes some of the very same points for which the expansion rate is negative when measured with the velocity field $`u_c^\mu `$. When this kind of ambiguity exists on the classical background there is little hope of finding a measure of spacetime expansion upon which all observers can agree when quantum perturbations are present.
A little thought leads to the realization that โ in attempting to identify a preferred velocity field from which to measure expansion in $`\mathrm{\Lambda }`$-driven inflation โ we have succumbed to the same confusion as the introductory physics student who seeks a preferred inertial frame. In the latter case, the laws of physics are phrased in terms of acceleration. In Galilean relativity, the co-ordinates of a given point in two different inertial frames of reference โ which are in relative motion with velocity $`\stackrel{}{V}`$ with respect to one another โ are related by:
$`\stackrel{}{r}`$ $`=`$ $`\stackrel{}{r}^{}+\stackrel{}{V}t,`$ (12)
$`t`$ $`=`$ $`t^{}.`$ (13)
The position and velocity at a given point are frame-dependent quantities while the acceleration is not.
So too in quantum gravity, it is really the local cosmological acceleration we wish to measure and not the expansion in some velocity field.
$``$ Acceleration in de Sitter: In general, the acceleration is obtained from the deviation equation of two infinitesimally close geodesics $`\chi ^\mu (\tau )`$ and $`\chi ^\mu (\tau )+\delta \chi ^\mu (\tau )`$:
$$\frac{D^2\delta \chi ^\mu (\tau )}{D\tau ^2}=R_{\nu \rho \sigma }^\mu \left[\chi (\tau )\right]\dot{\chi }^\nu (\tau )\delta \chi ^\rho (\tau )\dot{\chi }^\sigma (\tau ).$$
(14)
The curvature tensor for de Sitter spacetime equals:
$$deSitterR_{\beta \gamma \delta }^\alpha =H^2\left(\delta _\gamma ^\alpha g_{\beta \delta }\delta _\delta ^\alpha g_{\beta \gamma }\right).$$
(15)
We now consider two initially parallel, timelike geodesics with spacelike separation $`\mathrm{\Delta }^\mu `$. In synchronous gauge we have:
$`\chi ^\mu (\tau )`$ $`=`$ $`\tau \delta _0^\mu ,`$ (16)
$`\chi ^\mu (\tau )+\delta \chi ^\mu (\tau )`$ $`=`$ $`\tau \delta _0^\mu +\mathrm{\Delta }^\mu (0),`$ (17)
$`g_{\mu \nu }\dot{\chi }^\mu \mathrm{\Delta }^\nu =\mathrm{\hspace{0.33em}0}`$ , $`g_{\mu \nu }\dot{\chi }^\mu \dot{\chi }^\nu =1,`$ (18)
so that:
$$\frac{d\chi ^\mu (\tau )}{d\tau }\dot{\chi }^\mu =\delta _0^\mu .$$
(19)
For these geodesics, the deviation equation (14) takes the form:
$`{\displaystyle \frac{D^2\mathrm{\Delta }^\mu (\tau )}{D\tau ^2}}`$ $`=`$ $`H^2\mathrm{\Delta }^\mu g_{\rho \sigma }\dot{\chi }^\rho \dot{\chi }^\sigma +H^2\dot{\chi }^\mu g_{\rho \sigma }\dot{\chi }^\rho \mathrm{\Delta }^\sigma ,`$ (20)
$`=`$ $`H^2\mathrm{\Delta }^\mu ,`$ (21)
and a constant positive acceleration is manifest for any point on the manifold: the deviation between initially parallel, timelike and freely falling observers expands exponentially at all points in de Sitter spacetime. This is a frame invariant statement that characterizes the de Sitter geometry and, consequently, is free of the ambiguities that plague the velocity field.
$``$ Acceleration in general: It is convenient to use the freedom under general co-ordinate transformations to bring an arbitrary metric into synchronous gauge:
$$ds^2=dt^2+g_{ij}(t,\stackrel{}{x})dx^idx^j.$$
(22)
The geodesic deviation equation reduces to::
$`{\displaystyle \frac{D^2\mathrm{\Delta }^i(\tau )}{D\tau ^2}}`$ $`=`$ $`R_{0j0}^i\mathrm{\Delta }^j,`$ (23)
$`=`$ $`\left({\displaystyle \frac{1}{2}}g^{ik}\ddot{g}_{kj}+{\displaystyle \frac{1}{4}}g^{ik}g^{lm}\dot{g}_{kl}\dot{g}_{mj}\right)\mathrm{\Delta }^j,`$ (24)
The cosmological observation we are interested should not depend on the direction of the vector $`\mathrm{\Delta }^i`$, hence we contract into the vector by multiplying with $`g_{ij}\mathrm{\Delta }^j`$. Nor on the magnitude of the vector $`\mathrm{\Delta }^i`$, hence we scale the vector by dividing with its magnitude $`g_{rs}\mathrm{\Delta }^r\mathrm{\Delta }^s`$:
$$\frac{g_{ij}\mathrm{\Delta }^j}{g_{rs}\mathrm{\Delta }^r\mathrm{\Delta }^s}\frac{D^2\mathrm{\Delta }^i(\tau )}{D\tau ^2}=\frac{1}{g_{rs}\mathrm{\Delta }^r\mathrm{\Delta }^s}\left(\frac{1}{2}\ddot{g}_{ij}\frac{1}{4}g^{kl}\dot{g}_{ik}\dot{g}_{jl}\right)\mathrm{\Delta }^i\mathrm{\Delta }^j,$$
(25)
It is the right hand side of the geodesic deviation equation (14) which should provide the cosmological acceleration measurement $`\gamma `$ that the observer performs at event $`x`$:
$$\gamma (x)\frac{1}{g_{rs}\mathrm{\Delta }^r\mathrm{\Delta }^s}\left(\frac{1}{2}\ddot{g}_{ij}\frac{1}{4}g^{kl}\dot{g}_{ik}\dot{g}_{jl}\right)\mathrm{\Delta }^i\mathrm{\Delta }^j.$$
(26)
$``$ Acceleration for flat Robertson-Walker spacetimes: Let us restrict ourselves to the cosmologically interesting homogeneous, isotropic and spatially flat geometries:
$$FRWg_{ij}(t,\stackrel{}{x})a^2(t)\delta _{ij}.$$
(27)
Derivatives of the scale factor $`a(t)`$ give the velocity (Hubble) parameter $`H(t)`$ and the deceleration parameter $`q(t)`$:
$$H(t)\frac{\dot{a}}{a},q(t)\frac{a\ddot{a}}{\dot{a}^2}=1\frac{\dot{H}}{H^2}.$$
(28)
The relevant components of the curvature tensor are:
$$FRWR_{0j0}^i=\left(H^2+\dot{H}\right)\delta _j^i=qH^2\delta _j^i,$$
(29)
and, therefore, the cosmological acceleration observable $`\gamma `$ equals:
$$FRW\gamma (t)=qH^2=\frac{\ddot{a}}{a},$$
(30)
and it measures the fractional local cosmological acceleration.
$``$ Stochastic Acceleration: It is simple to evaluate the expectation value of the local acceleration $`\gamma `$ at one-loop order for the gravitational action (5) using stochastic techniques under the assumption of no back-reaction. First, we conformally re-scale the metric and express its time derivatives as functions of the re-scaled metric $`\stackrel{~}{g}_{ij}`$:
$`g_{ij}a^2\stackrel{~}{g}_{ij}`$ $``$ $`\dot{g}_{ij}=\mathrm{\hspace{0.33em}2}Ha^2\stackrel{~}{g}_{ij}+a^2\dot{\stackrel{~}{g}}_{ij},`$
$``$ $`\ddot{g}_{ij}=\left(2\dot{H}+\mathrm{\hspace{0.17em}4}H^2\right)a^2\stackrel{~}{g}_{ij}+\mathrm{\hspace{0.17em}4}Ha^2\dot{\stackrel{~}{g}}_{ij}+a^2\ddot{\stackrel{~}{g}}_{ij}.`$
We also define the Euclidean direction vector:
$$n^i\frac{\mathrm{\Delta }^i}{\sqrt{\mathrm{\Delta }^j\mathrm{\Delta }^j}}.$$
(32)
Our observable (26) becomes:
$$\gamma (x)=H^2+\dot{H}+\frac{1}{\stackrel{~}{g}_{rs}n^rn^s}\left(\frac{1}{2}\ddot{\stackrel{~}{g}}_{ij}+H\dot{\stackrel{~}{g}}_{ij}\frac{1}{4}\stackrel{~}{g}^{kl}\dot{\stackrel{~}{g}}_{ik}\dot{\stackrel{~}{g}}_{jl}\right)n^in^j.$$
(33)
Under the assumption of no back-reaction, the re-scaled spatial metric consists of the identity element plus just the transverse traceless field $`h_{ij}^{TT}`$ of linearized gravitons:
$$NoBackReaction\stackrel{~}{g}_{ij}=\delta _{ij}+\sqrt{32\pi G}h_{ij}^{TT},$$
(34)
and to linear order in $`G`$, the local acceleration operator is:
$`\gamma (x)`$ $`=`$ $`H^2+\dot{H}+[{\displaystyle \frac{1}{2}}\sqrt{32\pi G}\ddot{h}_{ij}^{TT}+H\sqrt{32\pi G}\dot{h}_{ij}^{TT}\mathrm{\hspace{0.17em}8}\pi G\dot{h}_{ik}^{TT}\dot{h}_{jk}^{TT}`$ (35)
$`(\mathrm{\hspace{0.17em}16}\pi G\ddot{h}_{ij}^{TT}+\mathrm{\hspace{0.17em}32}\pi GH\dot{h}_{ij}^{TT})h_{kl}^{TT}n^kn^l]n^in^j+O(G^{\frac{3}{2}}),`$
Taking the average over all directions and using:
$$n^in^j\frac{1}{3}\delta ^{ij},n^in^jn^kn^l\frac{1}{15}(\delta ^{ij}\delta ^{kl}+\delta ^{ik}\delta ^{jl}+\delta ^{il}\delta ^{jk}).$$
(36)
results in the following expression:
$`{\displaystyle d^2n\gamma (x)}`$ $`=`$ $`H^2+\dot{H}{\displaystyle \frac{8\pi G}{3}}\dot{h}_{ij}^{TT}\dot{h}_{ij}^{TT}{\displaystyle \frac{32\pi G}{15}}\ddot{h}_{ij}^{TT}h_{ij}^{TT}`$ (37)
$`{\displaystyle \frac{64\pi GH}{15}}\dot{h}_{ij}^{TT}h_{ij}^{TT}+O(G^{\frac{3}{2}}).`$
Before substituting the stochastic mode expansion for $`h_{ij}^{TT}`$ in (37) we must drop the terms where two derivatives act on the same field or upon different fields :
$$Stochasticd^2n\gamma (x)=H^2+\dot{H}\frac{64\pi GH}{15}\dot{h}_{ij}^{TT}h_{ij}^{TT}+O(G^{\frac{3}{2}})$$
(38)
Now from the expansion of the stochastic field in de Sitter spacetime : <sup>5</sup><sup>5</sup>5In (39), $`c.c.`$ denotes the conjugate expression.
$$h_{ij}^{TT}(t,\stackrel{}{x})=\frac{d^3k}{(2\pi )^3}\theta (Ha(t)k)\underset{\lambda }{}\{\frac{H}{\sqrt{2k^3}}e^{i\stackrel{}{k}\stackrel{}{x}}ฯต_{ij}(\stackrel{}{k},\lambda )\alpha _{\stackrel{}{k},\lambda }+c.c.\}$$
(39)
it is elementary to calculate its two-point function:
$`h_{ij}^{TT}(t,\stackrel{}{x})h_{ij}^{TT}(t^{},\stackrel{}{x})`$ $`=`$ $`{\displaystyle \frac{d^3k}{(2\pi )^3}\theta (Ha(t)k)\theta (Ha(t^{})k)}`$ (40)
$`\times {\displaystyle \frac{H^2}{2k^3}}{\displaystyle \underset{\lambda }{}}ฯต_{ij}(\stackrel{}{k},\lambda )ฯต_{kl}(\stackrel{}{k},\lambda ),`$
$`={\displaystyle \frac{2H^2}{16\pi ^3}}\mathrm{\hspace{0.17em}4}\pi {\displaystyle _H^{\mathrm{}}}{\displaystyle \frac{dk}{k}}\theta (Ha(t)k)\theta (Ha(t^{})k),`$
$`={\displaystyle \frac{H^2}{2\pi ^2}}\left\{\theta (tt^{})\mathrm{ln}\left[a(t^{})\right]+\theta (t^{}t)\mathrm{ln}\left[a(t)\right]\right\},`$
$`={\displaystyle \frac{H^3}{2\pi ^2}}\left[\theta (tt^{})t^{}+\theta (t^{}t)t\right].`$
The final answer emerges when we use (40) in expression (38):
$$deSitterd^2n\gamma (x)=H^2\left[\mathrm{\hspace{0.17em}1}\frac{32}{15\pi }GH^2+O(G^2)\right].$$
(41)
Even for GUT-scale inflation the correction is a very small constant effect. Because the stochastic formalism correctly captures only the leading infrared logarithms , our result (41) is consistent with zero change at one-loop order which is, in turn, consistent with the physics of back-reaction in quantum gravity on de Sitter spacetime.
Acknowledgements
It is a pleasure to acknowledge years of friendly and stimulating discussions on this topic with L. R. Abramo, R. H. Brandenberger, G. Geshnizjani, A. Guth, D. N. Page and W. G. Unruh. This work was partially supported by the European Social fund and National resources Y$`\mathrm{\Pi }`$E$`\mathrm{\Pi }\mathrm{\Theta }`$-PythagorasII-2103, by European Union grants FP-6-012679 and MRTN-CT-2004-512194, by NSF grant PHY-0244714, and by the Institute for Fundamental Theory at the University of Florida.
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# Quantum Electro and Chromodynamics treated by Thompsonโs heuristic approach
## I Introduction
There are a considerable number of problems in Science where fluctuations are present in all length scales, varying from microscopic to macroscopic wavelengths.
As examples, we can mention the problems of fully developed turbulent fluid flow, critical phenomena and elementary particle physics. The problem of non-classical reaction rates (diffusion limited chemical reactions) turns out also to be in this category.
As was pointed out by Wilson: โin quantum field theory, โelementaryโ particles like electrons, photons, protons and neutrons turn out to have composite internal structure on all sizes scales down to zero. At least this is the prediction of quantum field theoryโ.
The most largely employed strategy for dealing with problems involving many length scales is the โRenormalization - Group (RG) approachโ. The RG has been applied to treat the critical behavior of a system undergoing second order phase transition and has been shown to be a powerful method to obtain their critical indexes.
In Quantum Electrodynamics (QED), Gell-Mann and Low obtained a RG equation for electron charge $`e_\mu `$, being $`\mu `$ the energy scale, so that in the limit as $`\mu `$ goes to zero we obtain the classical electron charge e, and as $`\mu `$ goes to infinity we get the bare charge of the electron $`e_B`$.
The differential equation evaluated by Gell-Mann and Low obtains the โexperimentalโ charge $`e_\mu `$ of electron as a function of the energy, which corresponds to an interpolation between the classical and bare charges, namely: $`e<e_\mu <e_B`$.
In an alternative way to the RG approach, C. J. Thompson used a heuristic method (of the dimensions) as a means to obtaining the correlation length critical index ($`\nu `$), which governs the critical behavior of a system in the neighborhood of its critical point. Starting from Landau-Ginzburg-Wilson hamiltonian or free energy, he got a closed form relation for $`\nu (d)`$ , where d is the spatial dimension. It is argued that the critical behavior of this $`\mathrm{\Phi }^4`$-field theory is within the same class of universality as that of the Ising Model.
One of the present authors applied Thompsonโs method to study diffusion limited chemical reaction $`๐+๐\mathrm{๐}`$ (inert product). The results obtained in that work agree with the exact results of Peliti who renormalized term by term given by the interaction diagramms in the perturbation theory.
More recently, Nassif and Silva proposed an action to describe diffusion limited chemical reactions belonging to various classes of universality. This action was treated through Thompsonโs approach and could encompass the cases of reactions like $`๐+๐\mathrm{๐}`$ and $`๐+๐\mathrm{๐}`$ within the same formalism. Just at the upper critical dimensions of $`๐+๐\mathrm{๐}`$ ($`d_c=4`$) and $`๐+๐\mathrm{๐}`$ ($`d_c=2`$) reactions, the present authors found universal logarithmic corrections to the mean field behavior.
Thompsonโs renormalisation group method has been applied to obtain the correlation length critical exponent of the Random Field Ising Model by Aharony, Imry and Ma and by one of the present authors. His method was also used to evaluate the correlation length critical exponent of the N-vector Model. Yang - Lee Edge Singularity Critical Exponents has been also studied by this method.
The aim of this work is to apply Thompsonโs method to study some features of the QED and also QCD. As we will see, the evaluated $`QED_4`$ coupling $`\alpha `$ and the renormalized electron mass $`m`$ will exhibit logarithmic corrections on the energy scale. Here we would like to make some โalertโ as a form of caution: It is well known that the proper way of dealing with the problems of QED is by using the techniques of the Quantum Field Theory (QFT),by first treating the leading Feynman diagramms in the context of the perturbation theory. On the other hand, perturbation theory is not appropriate to deal with the infrared behaviour of the strong interaction described by the QCD.
QFT has also developed the apparatus for dealing with the divergences such as clever regularization schemes which leads to the renormalization of physical quantities to be measured. However the subject of this work is by far of more modest achievements.We have been exploring the various possibilities of the Thompsonโs method of dimensions (see refs ). As we can see , for instance, by considering these various possibilities of the method, we were able to obtain the universal logarithmic behavior for the coupling parameters of various models at their respective upper critical dimensions \[5,7,9-14\]. So now we would like to know how this method behaves when applied to $`QED_4`$ in order to obtain the logarithmic behavior on scale of coupling $`\alpha `$ (for $`d=4`$). The use of the present method here could be better justified taking into account the simplicity and universality of its application to study the behavior of coupling parameters of other models.
In section II, we start by first considering a heuristic prescription which characterizes Thompsonโs method. There we also introduce the QED lagrangian.
Third section is dedicated to some further elaboration of Thompsonโs method when applied to $`QED_4`$.
In the fourth section we study the equation describing the dependence both of the coupling constant $`\alpha `$ and the mass $`m`$ on the energy scale ($`\mu `$). Finally the next one will be dedicated to Quantum Chromodynamics (QCD) where we will evaluate the value of condesate of quarks given in its fundamental state and also the vacuum pressure at the boundary of the nucleon. Such results are motivated by the MIT-bag model , developed to describe the strong interaction inside the hadronic matter (nucleons).
The last section is dedicated to the conclusions and prospects.
## II $`QED`$ Lagrangian under the Thompsonโs method viewpoint
In this section, we start by writting the physical QED lagrangian, namely:
$$L=i\overline{\mathrm{\Psi }}\gamma ^\mu _\mu \mathrm{\Psi }m\overline{\mathrm{\Psi }}\mathrm{\Psi }\frac{1}{4}F^{\mu \nu }F_{\mu \nu }+ie\overline{\mathrm{\Psi }}\gamma ^\mu A_\mu \mathrm{\Psi }$$
(1)
with
$$F^{\mu \nu }=^\mu A^\nu ^\nu A^\mu and\overline{\mathrm{\Psi }}=\mathrm{\Psi }^{}\gamma ^{}$$
(2)
In(1) $`\mathrm{\Psi }`$ are fermion fields, $`e`$ and $`m`$ are respectively the electronโs charge and rest mass, $`A_\mu `$ is the four-vector electromagnetic potential and $`\gamma ^\mu `$ are the Diracโs matrices.
We assume that a heuristic approach used by Thompson to study critical phenomena can be applied to lagrangian (1). It states that:
โWhen we consider the integral of the Lagrangian (1) in a coherence volume $`l^d`$ in d-dimensions, the modulus of each integrated term of it is separately of the order of unityโ.
This method was firstly applied by Thompson to the Landau-Ginzburg-Wilson free energy or Hamiltonian, obtaining critical exponents within the same universality class of the Ising model.
But in fact, when we consider the integrals of each term in (1) as of the order of unity, we are really making a certain scaling dimensional analysis in each term of it. In doing this we have performed some scaling averages obtained separately from each integrated term of the lagrangian. In order to best justify the Thompsonโs prescription,we make the following reasoning: It is well known that $`QED_4`$ has a trivial fixed point ($`\alpha 0`$) at long- wavelength regimes ($`\lambda l\mathrm{}`$). Therefore if we want to describe the running coupling constant at low energy scales, it seems to be a good hypothesis to consider the behavior of the $`QED_4`$ lagrangian in the neighborhood of its trivial fixed point. This scaling invariance, which is also shared by a cooperative system at the critical point has been used by Thompson in order to obtain the correlation length critical index of the Ising-like systems,obtained close to its non-trivial fixed point. Indeed it is also well known that such a scaling invariance is in agreement with the canonical dimensions obtained here for the fields $`[\mathrm{\Psi }^2]_l(l^3)`$,$`[A_\mu ]_l(l^1)`$,and the coupling constant ($`[\alpha ]_ll^0constant`$) close to the fixed point ($`\alpha 0`$). So we conclude that the scaling prescription of Thompson, which comes from a dimensional analysis, works well in the neighborhood of a given fixed point.
We borrow Thompsonโs idea to apply it to $`QED_4`$, but in order to better do this, we make the following considerations:
As we consider the terms which appear in (1) to be quadratic forms like $`[\overline{\mathrm{\Psi }}\mathrm{\Psi }]`$ and $`[F^{\mu \nu }F_{\mu \nu }]`$, we are going to take Thompsonโs prescription into a more refined form, so that we can put the modulus of each integrated term exactly equal to the unity since these terms are taking in equal-footing. Besides this we will perform the integrations in the four -dimensional (4-D) space-time.
On the other hand, the idea of dimensional analysis is very common and when applied to evaluate the dimension of $`L`$ in QED leads to โ$`[L]=l^d=\mathrm{\Lambda }^d`$โ, in d-dimensions, where $`l`$ is the length and $`\mathrm{\Lambda }`$ is the momentum. By applying this prescription to each term of $`L`$, we obtain from the first one the dimension of the field $`\mathrm{\Psi }^2`$, that is simply $`[\mathrm{\Psi }^2]=\mathrm{\Lambda }^{d1}`$, which gives $`[\mathrm{\Psi }^2]=\mathrm{\Lambda }^3=l^3`$ for $`d=4`$.
In a similar way, from the third term of $`L`$ we get $`[A_\mu ^2]=\mathrm{\Lambda }^{d2}`$, being $`[A_\mu ^2]=\mathrm{\Lambda }^2=l^2`$ in the special case $`d=4`$.
Thompsonโs approach is based on a dimensional analysis in the energy scale (scaling reasoning) plus some additional heuristic considerations which lead to some mean values on the scale $`l`$ for the field $`\mathrm{\Psi },A_\mu `$, mass and charge (coupling $`\alpha `$).
Now applying Thompsonโs scaling assumption to the first term of (1) we have the following scaling integral:
$$\left|_{l^4}i[_\mu ][\overline{\mathrm{\Psi }}\mathrm{\Psi }]d^4x\right|=1.$$
(3)
We can observe that the dimension of ยด$`\gamma ^\mu _\mu `$\` ($`[\gamma ^\mu _\mu ]_l`$) which would appear in the integral is the same as $`[_\mu ]_l=l^1`$. This is because we are thinking only about a dimensional analysis in (3) for ยด$`\gamma ^\mu _\mu `$\`. So in this case we can naturally neglect the spinorial aspect of the field and just consider the ยดfirst derivative $`_\mu `$\`, which defines the fermions regarding to the scaling dimensional analysis: $`[_\mu ]_l=l^1`$.
On the other hand, when we are dealing with scalar fields, the second derivative picks up the bosonic behavior in a scaling dimensional analysis, i.e, $`[^\mu _\mu ]_l=[_\mu ^2]_l=l^2`$.
We would like to stress that, in his treatment of the critical phenomena, Thompson has considered integrals of the kind given by (3) in a more general case of d-dimensions. However, our interest here is restricted to the four-dimensional case, which is the most relevant if we take into account relativity theory, namely QED in (3+1) dimensions.
It is interesting to note that the integral (3) leads immediately to a kind of scaling dimensional analysis, where the dimensional value of certain quantity $`[\overline{\mathrm{\Psi }}\mathrm{\Psi }]`$ inside the integral is taken out of its integrand as a mean value in a coherent hyper-volume scale $`l^4`$. Then, from (3) we extract the following scaling behavior:
$`[\overline{\mathrm{\Psi }}\mathrm{\Psi }]_l[\overline{\mathrm{\Psi }}\mathrm{\Psi }]_l=[\mathrm{\Psi }^2]_ll^3`$, which also corresponds to a mean value of $`[\overline{\mathrm{\Psi }}\mathrm{\Psi }]`$ on scale $`l`$, where we have considered a 4-D hyper-cubic volume $`l^4`$ for (3), being $`[_\mu ]_l=l^1`$.
In order to apply Thompsonโs prescription to the second term of (1), we would consider the scaling $`\left|_{l^4}[m\overline{\mathrm{\Psi }}\mathrm{\Psi }]_xd^4x\right|=1`$. However, a close examination reveals that this procedure does not work quite well. Due to the coupling between the $`\mathrm{\Psi }`$ and $`A`$ fields, it is the mass increment that must be considered in the above relation. By putting $`\mathrm{\Delta }m=m(\mu )m_0`$, being $`\mu `$ the energy scale ($`\mu =l^1`$), $`\mathrm{\Delta }m`$ goes to zero as $`\mu 0`$ (or equivalently $`l\mathrm{}`$). After this consideration we can write:
$$\left|_{l^4}[\mathrm{\Delta }m]_x[\overline{\mathrm{\Psi }}\mathrm{\Psi }]_xd^4x\right|=1.$$
(4)
It is worth to emphasize that relation (4) was written by making the requirement that the quantities involved in it must satisfy a scaling relation. Relation (4) implies that:
$$[\mathrm{\Delta }m]_l[\overline{\mathrm{\Psi }}\mathrm{\Psi }]_ll^4=1,$$
(5)
or simply
$$[\mathrm{\Delta }m]_l[\overline{\mathrm{\Psi }}\mathrm{\Psi }]_ll^4=1.$$
(6)
As we know, $`[\overline{\mathrm{\Psi }}\mathrm{\Psi }]_l`$ goes as $`l^3`$, such that,from (6), we get that $`[\mathrm{\Delta }m]_l`$ scales as $`l^1`$, i.e, $`[\mathrm{\Delta }m]_ll^1`$.
To use Thompsonโs assumption in the third term of (1), it is better to think in terms of scaling behavior for the density of energy in the electromagnetic field ($`\rho `$), that is to say:
$$\frac{1}{8\pi }_{l^4}([E^2]_x+[B^2]_x)d^4x=1,$$
(7)
Relation (7) implies that $`[E^2]_l=[B^2]_ll^4`$. We know that $`\stackrel{}{B}=\times \stackrel{}{A}`$. So making a dimensional analysis for $`A`$, we obtain: $`[A^2]_l=[B^2]_ll^2l^4l^2=l^2`$.
Now let us consider the last term of the lagrangian (1). At a first sight, we could take a scaling integral for this term on a 4-D hyper-sphere($`l^4`$). Once the field $`A_\mu `$ in this term is a fluctuating field due to the fact that photons are emitted and absorbed by the electron, the mean value of $`A_\mu `$ on scale of suficiently long times would vanishes ($`A_\mu =0`$). So we need to consider something quadratic like $`A_\mu ^2`$ in order to avoid null mean value, or in other words, we need to evaluate a kind of second moment of this quantity. So based on this reasoning, as a means to extract a physical information (a quadratic coupling like $`e^2=\alpha `$) on the last term of (1), we judge necessary to perform a kind of second moment for the interaction term by considering an effective contribution for the action through a product of integrals, i.e, we must look for a product of integrals in a 4-dimensional space, which corresponds to an average of the square of the last term of (1) in an effective space of 8- dimensions. Making these considerations we firstly write
$$\left|i^2_{l^4}[_{l^4}(e[\overline{\mathrm{\Psi }}\mathrm{\Psi }]A_\mu )d^4x](e^{}[\overline{\mathrm{\Psi }}^{}\mathrm{\Psi }^{}]A_\mu ^{})d^4x^{}\right|=1,$$
(8)
where โ$``$โ is a dummy index.
We can also write (8) in the following way:
$$\left|i^2_{l^8}(e^2[\overline{\mathrm{\Psi }}\mathrm{\Psi }]_x^2[A_\mu ^2]_x)d^8x\right|=1,$$
(9)
which simply represents the following scaling: $`[(e[\overline{\mathrm{\Psi }}\mathrm{\Psi }]_l[A_\mu ]_l)l^4]^2=1`$.
## III Some further elaboration of Thompsonโs method applied to $`QED_4`$.
By reasons of spatial symmetry, let us now consider integral (3) evaluated in a volume of a 4 - D hyper-sphere, once we are interested in the isotropic 4-D space-time, being the scale of length $`l`$ the radius of this hyper-sphere.
The volume of a n-D hyper-sphere is given by $`V_n=S_n\frac{l^n}{n}`$ where $`S_n=\frac{2\pi ^{\frac{n}{2}}}{\mathrm{\Gamma }(\frac{n}{2})}`$. For 4-D, we have $`V_4=\frac{\pi ^2l^4}{2}`$, implying $`dV_4=2\pi ^2r^3dr`$, where $`r`$ is a radial variable.
The above considerations permit us to write integral (3) as being
$$2\pi ^2_{0_{(V_4)}}^l_r[\overline{\mathrm{\Psi }}\mathrm{\Psi }]_rr^3dr=[\overline{\mathrm{\Psi }}\mathrm{\Psi }]_l6\pi ^2_0^lr^2๐r=1.$$
(10)
Equation (10) implies that
$$[\overline{\mathrm{\Psi }}\mathrm{\Psi }]_l[\overline{\mathrm{\Psi }}\mathrm{\Psi }]_l=\frac{1}{2\pi ^2l^3},$$
(11)
where $`2\pi ^2l^3`$ is the magnitude of the surface of this 4-D hyper-sphere centered over the point charge $`e`$ whose โfieldโ (fermionic field) is given by the amplitude $`[\overline{\mathrm{\Psi }}\mathrm{\Psi }]_l`$ above.
$`[\overline{\mathrm{\Psi }}\mathrm{\Psi }]_l`$โ could be thought of as a mean of the squared fermionic field where the average is taken on a length scale $`l`$, being $`l=\mu ^1`$. Therefore this squared fermionic field has the dimension of $`\mu ^3`$ (the third power of energy). The โ$`2\pi ^2`$โ constant is a simple consequence of the spherical symmetry we have assumed for the problem.
Now let us evaluate the mass term given by integral (4) in the volume of a 4-D hyper-sphere of radius $`l`$. So we write
$$\left|2\pi ^2_0^l[\mathrm{\Delta }m]_r[\overline{\mathrm{\Psi }}\mathrm{\Psi }]_rr^3๐r\right|=1.$$
(12)
Relation (12) implies that
$$[\mathrm{\Delta }m]_l[\overline{\mathrm{\Psi }}\mathrm{\Psi }]_l\frac{\pi ^2l^4}{2}=1.$$
(13)
By putting (11) into (13), we get:
$$[\mathrm{\Delta }m]_l[\mathrm{\Delta }m]_l=4l^1.$$
(14)
Now let us consider the third term of lagrangian (1) and by performing the integral (7) in the volume of a 4-D hyper-sphere of radius $`l`$ ($`dV_4=2\pi ^2r^3dr`$),where the squared fields have the same scaling behavior $`[E^2]_l=[B^2]_ll^4`$, so we write:
$$\frac{\pi }{2}_0^l[E^2]_rr^3๐r=\frac{\pi }{2}_0^l[B^2]_rr^3๐r=1.$$
(15)
Relation (15) leads to
$$[E^2]_l=[B^2]_l=\frac{8}{\pi l^4}(l^4).$$
(16)
From the definition $`\stackrel{}{B}=\stackrel{}{}\times \stackrel{}{A}`$, we are led to the following scaling relation for $`[A^2]_l`$:
$$[B^2]_ll^2=[A^2]_l=\frac{8}{\pi l^2},$$
(17)
where in obtaining (17), we have used (16).
It is interesting to note that (17) is consistent with a potential of a static point charge, that is to say $`\mathrm{\Phi }\frac{1}{r}`$, which leads to $`[\mathrm{\Phi }^2]_l\frac{1}{l^2}`$, where $`\mathrm{\Phi }A_4`$ and $`A_\mu =(\stackrel{}{A},\mathrm{\Phi })`$. These considerations permit us to write (17) in a compact form:
$$[A_\mu ^2]_l=\frac{8}{\pi l^2}(l^2).$$
(18)
With respect to the fourth term of $`L`$ in (1), previous considerations had led to integral given by (9). Therefore let us now evaluate integral (9) in a volume of a 8-D hyper-sphere. Taking into account that $`d^8xdV_8=S_8r^7dr=\frac{\pi ^4}{3}r^7dr`$, we have:
$$\frac{\pi ^4}{3}_l[\alpha ]_r[\overline{\mathrm{\Psi }}\mathrm{\Psi }]_r^2[A_\mu ^2]_rr^7๐r=1,$$
(19)
where $`\alpha =e^2`$.
A first trying in order to evaluate (19) could be to write it as a product of averages, namely as a product of the quantities $`[\alpha ]_l,[\overline{\mathrm{\Psi }}\mathrm{\Psi }]^2_l,[A_\mu ^2]_l`$ and $`[V_8]_l`$. By considering that $`[\overline{\mathrm{\Psi }}\mathrm{\Psi }]^2_ll^6`$, $`[A_\mu ^2]_ll^2`$ and $`[V_8]_ll^8`$, we obtain that $`[\alpha ]_ll^6l^2l^81`$, which implies that $`[\alpha ]_l`$ is a constant, that is to say a quantity which does not exhibit a dependence on the scale of length $`l`$ (or energy $`\mu `$): $`[\alpha ]_ll^0`$ constant (scale invariance over the trivial fixed point).
We must consider that โ$`d=4`$โ corresponds to a kind of upper critical dimension for QED. In other words, below $`d=4`$ , fluctuations are very important to the problem, and above $`d=4`$, โmean fieldโ description is a good description to the problem. So $`d=4`$ represents a border-line dimension for QED and we must improve our approximations in order to โseeโ the logarithmic dependence of the coupling $`[\alpha ]_l`$ on the length scale $`l`$, or equivalently on the energy scale $`\mu =l^1`$. It must be stressed that a similar situation has been occurred when we treated diffusion limited chemical reactions through Thompsonโs method . As a means to improve the calculation of (19) let us take the quantities $`[\overline{\mathrm{\Psi }}\mathrm{\Psi }]_r^2`$ and $`[A_\mu ^2]_r`$ inside the integral with the same form as those evaluated in (11) and (18), but now displaying a dependence on the r-variable of scale. So by taking inside the integral (19) $`[\overline{\mathrm{\Psi }}\mathrm{\Psi }]_r=\frac{1}{2\pi ^2r^3}`$ and $`[A_\mu ^2]_r=\frac{8}{\pi r^2}`$, we can write:
$$\frac{\pi ^4}{3}_l[\alpha ]_r\left(\frac{1}{2\pi ^2r^3}\right)^2\left(\frac{8}{\pi r^2}\right)r^7๐r=1.$$
(20)
From (20),by putting the mean value of $`\alpha `$ on scale $`l`$ out of this integral, we have:
$$[\alpha ]_l\frac{2}{3\pi }_1^l\frac{dr}{r}=[\alpha ]_l\frac{2}{3\pi }ln(l)=1.$$
(21)
In evaluating (21), we have taken 1 as a lower cutoff on the scale $`l`$. Therefore (21) displays the logarithmic dependence for $`[\alpha ]_l`$ on the scale of length $`l`$ (or energy $`\mu =l^1`$).
## IV Evaluation of the dependence of charge and mass of the electron with the scale of energy
### A Obtaining $`\alpha (\mu )`$
In the quantum regime (vacuum polarization), the behavior of $`\alpha `$ is given by Eq.(21).
For the sake of simplicity in the notation, we write:
$$[\alpha ]_l[\alpha ]_l\alpha (l),$$
(22)
and by putting $`\mu =l^1`$ into (21), we get
$$\frac{2}{3\pi }ln(\mu )=\alpha ^1(\mu ).$$
(23)
Differentiating both sides of (23) with respect the $`\mu `$ variable, we obtain:
$$\mu \frac{d\alpha }{d\mu }=\frac{2}{3\pi }\alpha ^2.$$
(24)
Equation (24) coincides with that which is obtained by the R.G procedure, when $`QED_4`$ is treated through the perturbation theory at one loop level.We remember that, in performing the calculations by using the first prescription of Thompsonโs method, we firstly derived a coupling constant which does not depend on the energy scale. As a means to seek for a running coupling constant, we need to make a fine-tunning (see equation (20)) which recovers some fluctuations corrections which are present on scales of moderate energy. For the case of higher energy scales, we have already made an estimative for running coupling constant taking into account some additional assumptions for Thompsonโs prescription in order to get stronger fluctuations corrections.
The idea of fine-tunning was also used before in some previous works ,,\[12-14\], and it allowed us to obtain simple logarithmic corrections on energy scale just at the upper critical dimensions of those models\[5,7,12-14\].
From (24) we observe that we have obtained the coefficient $`\beta =\frac{2}{3\pi }\alpha ^2`$ ,.
Performing the integration of (24), by considering the limits $`\mu _0`$ and $`\mu `$ for the energy scales and their respectives couplings $`\alpha (\mu _0)`$ and $`\alpha (\mu )`$, we obtain:
$$\alpha (\mu )=\frac{\alpha (\mu _0)}{1\frac{2}{3\pi }\alpha (\mu _0)ln\left(\frac{\mu }{\mu _0}\right)}.$$
(25)
We observe that (25) diplays the so-called Landauโs singularity, namely a finite value of the energy scale $`\mu _L`$, where $`\alpha (\mu _L)\mathrm{}`$.
As it is well known, Landauโs singularity is a non-physical effect and reveals the fact that the running coupling constant solution given by (25) is not appropriate when the energy scale approaches $`\mu _L`$.
We are led to think that at higher energies, equation (24) and its solution (25) must be modified in order to be free of the Landauโs singularity. In the usual perturbative scheme of calculation this is accomplished by considering the theory beyond one loop level (two or more loops).
Now let us look at (25). We observe that $`lim_{\mu 0}\alpha (\mu )=0`$. But this result seems to be purely of academic interest.
Indeed even at low energy scales, the departure of the classical behavior for $`\alpha (\mu )`$ starts when $`\mu >m_0`$, where $`m_0`$ is the electron rest mass. This corresponds to assume that the effect of vacuum polarization in the shielding of electron charge becomes important when we approach the electron closest than its Compton wavelength $`l_0\lambda _c=m_0^1`$. Therefore, from an experimental point of view, we must look for (25) at the lower energy scale regime, but with $`\mu \mu _0m_0`$, that is, we consider the parameter of energy scale $`\mu `$ fixed on the electron mass $`m_0`$ as a scale of reference, where $`\alpha (\mu _0)\alpha (m_0)\frac{1}{137}`$. So for moderates energies we can make the expansion of (25), obtaining
$$\alpha (\mu )=\alpha _0\left[1+\frac{2}{3\pi }\alpha _0ln\left(\frac{\mu }{\mu _0}\right)\right],$$
(26)
where $`\alpha _0=\alpha (\mu _0)\alpha (m_0)\frac{1}{137}`$, and $`\mu _0m_0`$.
The result above (26) is well known in the literature. See reference.
### B Attainment of $`m(\mu )`$
As a means to evaluate $`\mathrm{\Delta }m(\mu )`$ let us compare (4) and (8) by considering the shift $`\mathrm{\Delta }e^2=\mathrm{\Delta }\alpha `$ in (8) because $`\mathrm{\Delta }\alpha (\mu )`$ must be directly proportional to the mass shift given in (4), that is, $`\mathrm{\Delta }m\mathrm{\Delta }\alpha `$ so that in the very lower energies limit we have $`\mathrm{\Delta }m\mathrm{\Delta }\alpha =\mathrm{\Delta }e^20`$. Thus in doing that we obtain
$$\left|_{V_4}(\mathrm{\Delta }m)[\overline{\mathrm{\Psi }}\mathrm{\Psi }]_xd^4x\right|=\left|i^2_{V_4}(_{V_4^{}}(\mathrm{\Delta }e^2)[\overline{\mathrm{\Psi }}^{}\mathrm{\Psi }^{}]_x^{}[A_\mu ^2]_x^{}d^4x^{})[\overline{\mathrm{\Psi }}\mathrm{\Psi }]_xd^4x\right|=1,$$
(27)
where we consider the shift $`\mathrm{\Delta }m=mm_0`$ and $`\mathrm{\Delta }\alpha =\alpha \alpha _0`$, being $`m_0`$ the electron rest mass and $`\alpha _0\frac{1}{137}`$ a constant measured on energy scale of electron rest mass.
Relation (27) implies that
$$\mathrm{\Delta }m=\left|i^2_{V_4}(\mathrm{\Delta }e^2)[\overline{\mathrm{\Psi }}\mathrm{\Psi }]_x[A_\mu ^2]_xd^4x\right|,$$
(28)
where the index โ$``$โ is dummy.
By putting $`d^4xdV_4=2\pi ^2r^3dr`$, $`[\overline{\mathrm{\Psi }}\mathrm{\Psi }]_x[\overline{\mathrm{\Psi }}\mathrm{\Psi }]_r=\frac{1}{2\pi ^2r^3}`$ and $`[A_\mu ^2]_x[A_\mu ^2]_r=\frac{8}{\pi r^2}`$ (see (11) and (18)) into (28), we get
$$\mathrm{\Delta }m=\frac{8}{\pi }[\mathrm{\Delta }\alpha ]_l\left|_l\frac{1}{r^2}๐r\right|.$$
(29)
Now let us take the notation $`[\mathrm{\Delta }\alpha ]_l\mathrm{\Delta }\alpha `$ to represent a mean charge shift measured on the scale $`l\mu ^1`$. Thus performing the integration (29) between the limits $`\mu =0`$ ($`l=\mathrm{}`$) and $`\alpha _0l_0\alpha _0\lambda _c10^{14}m`$, which is equivalent to the vacuum polarization regime due to the small value of length scale $`\alpha _0ล_0`$ ($`<<\lambda _c`$,i.e, the Compton wavelength), we obtain
$$\mathrm{\Delta }m=\frac{8}{\pi }\mathrm{\Delta }\alpha \frac{1}{\alpha _0l_0}=\frac{8}{\pi }\frac{\mathrm{\Delta }\alpha }{\alpha _0}m_0$$
(30)
where $`m_0=l_0^1\lambda _c^1`$. Indeed we verify such a proportionality โ$`\mathrm{\Delta }m\mathrm{\Delta }\alpha `$โ (30) mentioned before, which leads to
$$m=m_0+\frac{8}{\pi }\frac{\mathrm{\Delta }\alpha }{\alpha _0}m_0,$$
(31)
where $`\mathrm{\Delta }m=mm_o`$.
Finally by substituting $`\mathrm{\Delta }\alpha (\mu )`$ obtained from (26) into (31) we get
$$m=m_0\left[1+\frac{16}{3\pi ^2}\alpha _0ln\left(\frac{\mu }{\mu _0}\right)\right].$$
(32)
We notice that the above relation is comparable to the result for $`m(\mu )`$ of the literature as quoted by Nottale, Weinberg and Weisskopf . So, in spite of all results we have just obtained are well-known in the literature, we must stress again that the novelty here consists in obtainning another and a new way to deal with some problems of QED by applying Thompsonโs method of scales.
It is also interesting to verify that the quantum correction to the electron mass could be obtained in a way which is consistent with (32) through the following reasoning: Let us consider the energy stored in the eletric field, namely $`U_{el}=_{V_3}E^2๐V_3`$ , where the integration is performed in a 3-D volume. Now let us write the eletric field $`E`$ as its classical value $`E_0`$ plus a correction $`\mathrm{\Delta }E`$ due to the quantum fluctuations. We assume that quantum fluctuations only affects the energy $`U_{el}`$ through the squared contribution in $`\mathrm{\Delta }E`$, once the linear term in $`\mathrm{\Delta }E`$ averages out to zero. So we have $`\overline{E^2}=\overline{E_0^2}+\overline{\mathrm{\Delta }E^2}`$ , where the bars means averaging over a suficiently long time in the scale of the fluctuations. Therefore, as we are mainly interested in the quantum process namely the absortion and emission of virtual photons, we can write
$$\mathrm{\Delta }E_{rms}=[\overline{(\mathrm{\Delta }E^2)}]^{\frac{1}{2}},$$
(33)
where the index $`rms`$ means root mean square.
It is natural to think that $`\mathrm{\Delta }E_{rms}`$ will be different from zero only in the presence of the fermionic field, and this leads us to propose the relation
$$\mathrm{\Delta }E_{rms}^2=\xi ^2\mathrm{\Psi }_{rms}^2,$$
(34)
where we have considered $`\mathrm{\Psi }_{rms}^2=[\overline{\mathrm{\Psi }}\mathrm{\Psi }]_r=\frac{1}{2\pi ^2r^3}`$(see (11)),that is to say ,$`\mathrm{\Psi }_{rms}^2`$ corresponds to the mean squared fermionic field in the variable of scale-$`r`$, and $`\xi `$ is a constant.
It seems that the intuitive reasoning given in (34) is consistent with gauge invariance of the theory. That is, in order to take into account the quantum fluctuations contribution to the electric field, we must have necessarily the coupling between the electric and the fermionic fields.
Equations (34) and (11) imply that
$$\mathrm{\Delta }E_{rms}\frac{1}{r^{\frac{3}{2}}},$$
(35)
which must be compared with the inverse square Gauss law of the classical contribution. At this point we would like to notice that a dependence of the quantum fluctuations of the eletric field on the scale of length as that we obtained in (35) was proposed by Weisskopf.
Now performing the integration of $`\mathrm{\Delta }E_{rms}^2`$ in a 3-D volume and taking as integration limits the variable r $`(r\lambda _c)`$ and $`\lambda _c`$ the Compton wavelength, we obtain
$$(\mathrm{\Delta }m)c^2=mc^2m_0c^2=_r^{\lambda _c}(\frac{\xi ^2}{2\pi ^2r^3})4\pi r^2๐r.$$
(36)
The reason to consider $`\lambda _c`$ as long wavelength cutoff is that quantum fluctuations does not contribute very much to the electromagnetic mass of the electron above this value.
Relation (36) implies that
$$mc^2=m_0c^2+Cln(\frac{\lambda _c}{r}),$$
(37)
where $`C`$ is a constant.
It is worth to emphasize that (37) is consistent with the results we have obtained in (32) and that by Weisskopf,if we fix $`C\frac{3}{2\pi }m_0c^2\alpha _0`$, being $`\frac{\lambda _c}{r}\frac{\mu }{\mu _0}`$.
## V Thompsonโs method, QCD and MIT-bag model
Quantum Chromodynamics (QCD), the modern theory of the strong interactions is a non-Abelian Field Theory. In 1973, Gross and Wilczek and independently Politzer have shown that certain classes of non-Abelian fields theories exhibit asymptotic freedom, a necessary condition for a theory which could describe strong interactions. These seminal papers open the route to the birth of the QCD.
In a not very accurated picture, QCD can be considered as an expanded version of QED. In QCD we have also six fermionic fields representing the various quark flavors, in contraposition to a simgle fermionic field of the QED. Besides the asymptotic freedom exhibit at the ultraviolet limit, a theory of the strong interactions must also display quark confinement at the infrared limit.
Whereas in QED there is just one kind of charge, QCD has three kinds of charge, labeled byโcolorโ (red, green and blue). The color charges are conserved in all physical process. There are also photon-like massless particles, called color gluons, that respond in appropriate ways to the presence of color charge. This mechanism is very similar to the ways photons respond to electric charge in QED.
In QCD, quarks are particles that carry color charge. As we already know, there are six different kinds of quarks, called โflavorsโ, denoted by $`u`$ (up), $`d`$ (down); $`c`$ (charmed), $`s`$ (strange); $`b`$ (botton) and $`t`$ (top). Of these, only $`u`$ and $`d`$ quarks play a significant role in the structure of ordinary matter. They carry fractional electric charge, i.e, $`+\frac{2}{3}e`$ for $`u`$, $`c`$ and $`t`$ quarks, and $`\frac{1}{3}e`$ for $`d`$, $`s`$ and $`b`$ quarks, in addition to their color charge.
In a similar way to QED-Lagrangian, let us write the QCD-Lagrangian density, namely:
$$L=\mathrm{\Sigma }_j\overline{q}_j(i\gamma _\mu D^\mu m_j)q_j\frac{1}{4}G_{\mu \nu }^aG_a^{\mu \nu },$$
(38)
where $`D^\mu =^\mu +\frac{1}{2}ig\lambda _aA_a^\mu `$ , and $`G_a^{\mu \nu }=^\mu A_a^\nu ^\nu A_a^\mu gf_{abc}A_b^\mu A_c^\nu `$
In (38) above, $`m_j`$ and $`q_j`$ are the mass and quantum field of the quark of $`j^{th}`$ flavor, and $`A`$ is the gluon field, being $`\mu `$ and $`\nu `$ the space- time indices. $`a`$, $`b`$ and $`c`$ are color indices. The numerical coefficients $`f`$ (structure constants) and $`\lambda _a`$ guarantee $`SU(3)`$ color symmetry. $`g`$ is the coupling constant.
By applying Thompsonโs assumption to the Knetics term from (38), we obtain a similar result as given before in QED (see equation (11)), namely
$$[\overline{q}_jq_j]_l[\overline{q}_jq_j]_l=\frac{1}{2\pi ^2l^3}.$$
(39)
But now, the amplitude above must be interpreted as a quark condensate scaling, and we took the negative signal in order to be consistent with the description of a bound state.
On the other hand, let us remember that we have obtained the mass shift $`\mathrm{\Delta }m`$ of the electron in QED (see equation (29)). This is due to the interaction with electromagnetic field (photon). Following this same reasoning for QCD, where we have quarks (fermions) interacting with gluon fields (bosons), thus according to MIT-bag model idea, now let us think of a special behavior for the running coupling constant ($`\alpha _{running}`$) which obeys a step-like function specified by considering the boundary condition at the surface of the bag, namely:
$$\alpha _{runing}=\{\begin{array}{cc}1,\hfill & \frac{1}{m_0}<r<\mathrm{}\hfill \\ 0,\hfill & r<\frac{1}{m_0}(=r_0),\hfill \end{array}$$
(40)
where $`m_0`$ is a reference mass above which the running coupling constant vanishes, i.e, for a radius below nucleon radius ($`r<r_0`$) or $`m>m_0`$,the quarks have free motion inside the bag (the nucleons); so for $`r>r_0`$ we consider $`\alpha =1`$.
Finally, taking into account $`\alpha _{running}`$ in (40) inside the integral for $`\mathrm{\Delta }m`$ (equation (29)), and making the integration between the limits $`r=\mathrm{}`$ and $`r=r_0=\frac{1}{m_0}`$, we write
$$\mathrm{\Delta }m=\frac{8}{\pi }\left|_{\mathrm{}}^{\frac{1}{m_0}}\alpha _{running}\frac{1}{r^2}๐r\right|,$$
(41)
where we obtain
$$\mathrm{\Delta }m=\frac{8}{\pi }m_0.$$
(42)
For the special case of strong interaction ($`\alpha =1`$), we consider that the shift of mass, which comes from the strong interaction effects inside the nucleon, is pratically responsable for almost all the mass of nucleon. So due to this fact, we can make the following approximation: $`\mathrm{\Delta }mm_{nucleon}=m_n`$. Thus by introducing it inside (42), we get
$$m_n=\frac{8}{\pi }m_0.$$
(43)
Now using $`h=c=1`$, and also by considering $`m_0`$ as a zero-point energy inside (43), i.e, we have $`m_0=\frac{1}{2}\nu _0=\frac{1}{2l}`$ inside (43), thus we obtain:
$$m_n=3m_q=\frac{4}{\pi l},$$
(44)
where $`m_q`$ represents the constituent mass of quark ($`m_q=\frac{1}{3}m_n`$).
From (44) above, we get $`\frac{1}{l}=\frac{3\pi }{4}m_q`$. So by putting this result inside the scaling relation for the quark condensed given in (39), we obtain
$$[\overline{q}_jq_j]_l[\overline{q}_jq_j]_l=\frac{27\pi }{128}m_q^3.$$
(45)
Having $`m_n939MeV`$, which implies that $`m_q=\frac{1}{3}m_n313MeV`$, we finally get $`[\overline{q}q](\frac{m_q}{1.147}MeV)^3(273MeV)^3`$.
The value of the quark condensate evaluated above must be compared with other theoretical values which go from $`(265MeV)^3`$ to $`(340MeV)^3`$ (see table (2) in the paper by Mota $`etal.`$) and also with recent experimental result of $`[(296\pm 25)MeV]^3`$.
According to MIT bag-model, the simplest shape for a bag is naturally a sphere (spherical bags), i.e, $`R(\theta ,\varphi )=R=constant`$; $`๐ง(\theta ,\varphi )=๐_r`$ is the unitary normal vector to the spherical surface. $`R`$ is the bag radius ($`Rr_0`$). Another boundary condition is obtained by demanding that the pressure of the quarks on the bag surface be constant and must equal a constant exterior vacuum pressure, i.e, we have $`B=\frac{1}{2}๐ง.\mathrm{\Sigma }_q\overline{q}q]_{R=R(\theta ,\varphi )}`$.
On the other hand, as we already know, the condensed $`[\overline{q}q]`$ was obtained by Thompsonโs scaling reasoning (equation 39). So inserting (39) into the formula above from MIT-bag model, and by considering $`๐ง.=\frac{}{r}`$, thus we get
$$B=\frac{1}{2}\frac{}{r}[\overline{q}q]_r=\frac{1}{2}\frac{d}{dr}(\frac{1}{2\pi ^2r^3})=\frac{3}{4\pi ^2r^4}.$$
(46)
Having $`hc=2\pi \mathrm{}c=1`$, then by introducing this information into (46) above, we finally write
$$B=\frac{6\mathrm{}c}{4\pi r^4}.$$
(47)
The above result, which comes from MIT-bag model in addition to the scaling result for quark condensate in spirit of Thompsonโs approach (equation (39)), coincides with that obtained in another work, dealing with the MIT-bag model.
Since nucleon has a radius in order of $`10^{15}m`$, thus from (47) we obtain $`B10^{29}atm`$, which also represents an external vacuum pressure over the nucleon. Such a pressure must equal the pressure inside the bag (nucleon) due to the quarks.
## VI Conclusions and prospects
In this paper, Thompsonโs heuristic method which could be considered as a simple alternative way to the RG calculations was applied to study $`QED_4`$.This was done by treating each term of the QED lagrangian in equal footing, through dimensional analysis on the scale of length (or equivalently on the momentum-energy scale). If we analyse the scaling behavior of certain objects such that the mean squared fermionic field ($`[\overline{\mathrm{\Psi }}\mathrm{\Psi }]_l`$), the dimension of the squared vector potencial $`[A_\mu ^2]_l`$, the โexcessโ of mass $`[\mathrm{\Delta }m]_l`$, and the โexcessโ of charge $`[\mathrm{\Delta }\alpha ]_l`$, with all these quantities evaluated at the scale of length $`l`$, we observe that it is possible to organize these objects within a hierarchical structure, thinking in terms of topological grounds. In this way, the mean squared fermionic field $`[\overline{\mathrm{\Psi }}\mathrm{\Psi }]_l=(2\pi ^2l^3)^1`$ decreases as a โsurfaceโ 3-D of a hyper sphere 4-D of radius $`l`$, being this โsurfaceโ immersed in the 4-D space-time.
The next object in this hierarchy corresponds to the dimension of the squared vector potential. It is given by $`[A_\mu ^2]_l=8(\pi l^2)^1`$, exhibiting a inverse square law on the scale of length $`l`$. This represents a 2-D structure also immersed in the 4-D space-time. We could think that, for this object, the degree of freedom have reduced by a unity. The โexcessโ of mass $`[\mathrm{\Delta }m]_l=4l^1`$ can be thougth of as a 1-D structure immersed again in a 4-D space-time.
Finally the โexcessโ of charge (coupling) $`[\mathrm{\Delta }\alpha ]_l`$ behaves in the zero-th order as scaling independent, namely $`\alpha `$ goes as $`l^0`$ at zero order in the calculations, and it can be considered as a 0-D structure immersed in a 4-D space-time. In short, we have the โspreadingโ of the squared fermionic field in a 3D-volume , the squared vector potential in a 2D-surface , the mass in a 1D-line and the charge in a point (0-D), relating these objects of QED to a hierarchical ordering in the topology of a 4-D space-time.
However when we improve our calculations, the charge (running coupling โconstantโ) passes to exhibit a logarithmic dependence on the scale of length. This could be considered as an intermediate regime between a point ($`l^0`$ constant) and a line ($`l^1`$). We interpret this as the charge accquiring a fractal character in this topological structure of the space-time, due to the influence of the quantum fluctuations introduced by the vacuum polarization, in such a way that we have $`\alpha (l)[ln(l)]^1=[l^0ln(l)]^1`$. These quantum fluctuations will also โmodulateโ the behavior of the โexcessโ of mass, namely $`\mathrm{\Delta }m(l)[l^1ln(l)]^1`$.
The fractal character of a quantum path was considered by Nottale on analysing the QED. He showed that due to the vacuum polarization the self-energy diagramms of the QED display a fractal character.
One merit of Thompsonโs approach is that it displays the scaling behavior of the physical magnitudes of the problem, and as a consequence, it allows us for instance to pick up the fractal structure of $`\alpha (\mu )`$ and $`m(\mu )`$ (logarithmic dependence on scale).
One of the possibilities of the Thompsonโs method is to use it as a means to evaluate the running coupling constant of Quantum Chromodynamics (QCD). This matter will be treated elsewhere.
Finally, since Thompsonโs method is essentially a scaling approach, we can also apply it to study the growth of polymer chains in an alternative way to the R.G scheme. This matter will be also treated elsewhere.
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# Exact Results for Three-Body Correlations in a Degenerate One-Dimensional Bose Gas
## Abstract
Motivated by recent experiments we derive an exact expression for the correlation function entering the three-body recombination rate for a one-dimensional gas of interacting bosons. The answer, given in terms of two thermodynamic parameters of the Lieb-Liniger model, is valid for all values of the dimensionless coupling $`\gamma `$ and contains the previously known results for the Bogoliubov and Tonks-Girardeau regimes as limiting cases. We also investigate finite-size effects by calculating the correlation function for small systems of 3, 4, 5 and 6 particles.
Quantum fluctuations are well known to have a profound influence on the physics of one-dimensional (1D) systems. Due to quantum fluctuations, a 1D gas of weakly interacting bosons does not exhibit true long-range order and hence does not undergo Bose-Einstein condensation at any temperature. With increasing repulsion between the bosons, the gas exhibits a smooth crossover to the strong coupling regime, where the correlations in the positions of the particles endow it with fermion-like thermodynamic properties. A strong repulsion between the bosons also has a marked effect on the local fluctuation properties, leading, for example, to a fermion-like suppression of the fluctuations of particle density. In this paper we shall focus on a three-body correlation function, which is directly related to the lifetime of a Bose gas unstable with respect to three-body recombination processes. It has been studied in a recent experiment, where measurements of the three-body recombination rate for trapped atoms were performed Tolra .
The idea of using three-body recombination for a measurement of local correlations was originally put forward by Kagan et al. KSS-85 . These authors showed that the low-temperature recombination rate for a 3D Bose gas is proportional to the local three-body correlation function $`g_3`$ defined by the ground-state expectation value
$$g_3=:\widehat{n}^3(๐ซ):\widehat{\mathrm{\Psi }}^{}(๐ซ)^3\widehat{\mathrm{\Psi }}(๐ซ)^3.$$
(1)
Here $`\widehat{n}(๐ซ)=\widehat{\mathrm{\Psi }}^{}(๐ซ)\widehat{\mathrm{\Psi }}(๐ซ)`$ is the particle density operator, $`\widehat{\mathrm{\Psi }}^{}`$ ($`\widehat{\mathrm{\Psi }}`$) denote boson creation (annihilation) operators and the symbol $`::`$ stands for normal ordering. It was noted that, due to the factor-of-six suppression of $`g_3`$ in a Bose-Einstein condensate at zero temperature as compared to the non-condensed Bose gas, three-body recombination can be used as a diagnostic tool for distinguishing between the condensed and the non-condensed phases. This was later confirmed experimentally Burt .
In the experiment of Ref. Tolra the recombination process was used for the investigation of three-body correlations in a 1D Bose gas. A magnetically trapped Bose-Einstein condensate of <sup>87</sup>Rb atoms was loaded into a deep 2D optical lattice, by which the condensate was divided into an array of independent 1D systems. In the case of <sup>87</sup>Rb two-body losses are very small, and it was possible to model the decay in time of the number of trapped atoms in terms of one-body and three-body processes only. The rate of change of the total number of atoms $`N`$ was written as
$$\frac{dN}{dt}=K_1N๐๐ซK_3^{1\mathrm{D}}n_{3\mathrm{D}}^3,$$
(2)
where $`K_1`$ and $`K_3^{1\mathrm{D}}`$ are the rate coefficients for one- and three-body losses. By measuring the decrease of the number of atoms as a function of time it was found that the rate coefficient $`K_3^{1\mathrm{D}}`$ in the 1D gas was reduced considerably, by a factor of about 7, compared to its value $`K_3^{3\mathrm{D}}`$ for a 3D condensate. The measurements were carried out for a particular value of the dimensionless coupling constant $`\gamma `$ ($`0.5`$) characterizing the strength of the 1D correlations. Given that $`K_3`$ is proportional to the three-body correlation function $`g_3`$, the observed reduction was interpreted as an effect of the reduced dimensionality on $`g_3`$.
In the following we shall therefore focus on calculating $`g_3`$ for a 1D Bose gas. Our treatment is based on the Lieb-Liniger (LL) model Lieb for which the dimensionless coupling parameter $`\gamma `$ is given in Eq. (5) below. In the thermodynamic limit, where the particle number $`N`$ and the size of the system $`L`$ tend to infinity while the density $`n`$ remains constant,
$$n=N/L=\mathrm{const},N,L\mathrm{},$$
(3)
an expression for $`g_3`$ has so far been obtained only for small and large values of $`\gamma `$ Gangardt . Here we report an explicit expression for $`g_3,`$ valid for all $`\gamma .`$ Its derivation employs an integrable lattice regularization of the LL model BBP-93 ; BIK-98 together with conformal field theory. The calculations are quite involved and will be presented elsewhere CSZ . The answer is given by Eq. (17) in terms of two thermodynamic parameters of the LL model: the second- and fourth-order moment, Eq. (16), of the quasi-momentum distribution function. The quasi-momentum distribution function is the solution of the Lieb equation (14), which is a linear integral equation.
Within the LL model the Hamiltonian for $`N`$ identical bosons of mass $`m`$ is taken to be
$$H=\frac{\mathrm{}^2}{2m}\left[\underset{i=1}{\overset{N}{}}\frac{^2}{x_i^2}+2c\underset{1i<jN}{}\delta (x_ix_j)\right].$$
(4)
The interaction constant $`c0`$ has the dimension of inverse length. The dimensionless coupling strength $`\gamma `$ is given by
$$\gamma =c/n.$$
(5)
The LL model allows for an exact determination of the eigenfunctions and eigenvalues of $`H`$. In the $`\gamma \mathrm{}`$ limit of the model, known as the Tonks-Girardeau (TG) gas Tonks , the eigenvalues are the same as that of a free fermion gas.
While the eigenfunctions and eigenvalues are known exactly for the LL model, its correlation functions in general and $`g_3`$ in particular have been much less investigated. The reason for this is that the eigenfunctions are very complicated for a general value of $`N`$. We exhibit these below and calculate $`g_3`$ analytically for the simplest possible case $`N=3.`$
The eigenfunctions $`\mathrm{\Psi }(x_1,x_2,\mathrm{},x_N)`$ of the Hamiltonian (4) must be symmetric under any interchange of coordinates. Since the particles interact via a delta-function potential, the derivatives of $`\mathrm{\Psi }`$ jumps when two particles approach each other, while when all $`x_j`$ are different, $`\mathrm{\Psi }`$ satisfies the free-particle Schrรถdinger equation. The explicit expression for $`\mathrm{\Psi }`$ was obtained in Ref. Lieb and given in Ref. KBI-93 in the following form
$$\begin{array}{c}\mathrm{\Psi }(x_1,x_2,\mathrm{},x_N)=C\underset{P}{}(1)^{[P]}\mathrm{exp}\left\{i\underset{j=1}{\overset{N}{}}k_{P_j}x_j\right\}\hfill \\ \hfill \times \underset{j>l}{}[k_{P_j}k_{P_l}ic\mathrm{sgn}(x_jx_l)],\end{array}$$
(6)
where $`P`$ is a permutation of $`N`$ numbers and $`[P]`$ is the parity of the permutation. The possible values of the quasi-momenta $`k_j`$ are determined by the boundary conditions imposed on the system. For periodic boundary conditions (that is, for $`N`$ particles placed on a ring of circumference $`L`$) $`k_j`$ are solutions to the following system of coupled nonlinear equations, called the Bethe equations,
$$\mathrm{exp}\left\{ik_jL\right\}=\underset{lj}{\overset{N}{}}\frac{k_jk_l+ic}{k_jk_lic},j=1,\mathrm{},N.$$
(7)
The calculation of the normalization constant $`C`$ in Eq. (6) is a nontrivial task. For the periodic boundary conditions a closed expression for $`C`$ was suggested by Gaudin (see Gaudin-83 and references therein) and proved in Ref. Korepin-82 ; a detailed discussion can be found in Ref. KBI-93 . One has now all the ingredients for calculating $`g_3g_3(\gamma ,N)`$ for the 1D system. Written in the first-quantized form, $`g_3(\gamma ,N)`$ is
$$g_3(\gamma ,N)=\frac{N!}{(N3)!}_0^L๐X|\mathrm{\Psi }(0,0,0,x_4,\mathrm{},x_N)|^2,$$
(8)
where $`dX=dx_4\mathrm{}dx_N`$ and $`\mathrm{\Psi }`$ is given by Eq. (6). Note that the average $`\mathrm{}`$ in Eq. (1) is taken over the ground state of the system; equivalently, it is the ground-state wave function which enters Eq. (8). This is ensured by selecting the proper set of the quasi-momenta $`k_j`$ among all the solutions of the Bethe equations (7).
All quasi-momenta $`k_j`$ are equal to zero for the non-interacting system, $`\gamma =0,`$ and the wave function (6) is uniform in space, $`\mathrm{\Psi }=(1/\sqrt{L})^N`$ , which implies that
$$g_3(0,N)=n^3N(N1)(N2)/N^3.$$
(9)
We now use the wave function (6) to get $`g_3(\gamma ,N)`$ for $`N=3`$ analytically, and for $`N=4,5`$ and $`6`$ numerically. The quasi-momenta in the ground state of the system with $`N=3`$ obey $`k_1=k_3`$ and $`k_2=0.`$ It is convenient to write the Bethe equation for $`k_1`$ in the form
$$\lambda =2\pi 2\mathrm{arctan}(\lambda /3\gamma )2\mathrm{arctan}(2\lambda /3\gamma ),$$
(10)
where we have introduced $`\lambda k_1L`$ and the principal branch for arctangent has been chosen: $`\pi /2\mathrm{arctan}(x)\pi /2.`$ The solution to this transcendental equation grows monotonically from $`\lambda 3\sqrt{\gamma }`$ for $`\gamma 1`$ to $`\lambda 2\pi 4\pi /\gamma `$ for $`\gamma 1`$. The resulting answer for $`g_3(\gamma ,3)=6|\mathrm{\Psi }(0,0,0)|^2`$ is
$$g_3(\gamma ,3)=\frac{16}{3}\frac{A^2B\lambda ^6}{\gamma ^3L^3[1+4(3A^2+2A+B+6AB)]},$$
(11)
where $`A=3\gamma /(9\gamma ^2+\lambda ^2)`$ and $`B=3\gamma /(9\gamma ^2+4\lambda ^2).`$ The ratio $`g_3(\gamma ,3)/g_3(0,3)`$ is plotted in Fig. 1.
It decreases monotonously with increasing $`\gamma ,`$ from $`129\gamma /24`$ for $`\gamma 1`$ to $`512\pi ^6/243\gamma ^6`$ for $`\gamma 1`$. In Fig. 1 we also plot $`g_3(\gamma ,N)`$ calculated for $`N=4`$, $`5`$ and $`6.`$ Note that the magnitude of the slope at the origin ($`\gamma `$=0) increases with increasing number of particles, and when $`N\mathrm{}`$ it approaches infinity, in agreement with the Bogoliubov result given in Ref. Gangardt :
$$g_3(\gamma ,\mathrm{})/n^316\sqrt{\gamma }/\pi ,\gamma 0.$$
(12)
We have obtained an analytic expression for $`g_3`$ in the Tonks-Girardeau limit, $`\gamma \mathrm{},`$ for arbitrary values of $`N`$ by generalizing the result Eq. (16) of Ref. Gangardt , valid in the thermodynamic limit, to the case of a finite number of particles. The result is
$$\frac{g_3(\gamma ,N)}{n^3}\frac{16\pi ^6}{15\gamma ^6}\frac{(N^21)^2(N^24)}{N^6},\gamma \mathrm{}.$$
(13)
We now present our main result: the exact expression for the three-body correlation function $`g_3(\gamma ,\mathrm{})`$ in the thermodynamic limit (3) of the Lieb-Liniger model (4). The system of Bethe equations (7) reduces in this limit to the linear integral equation called the Lieb equation:
$$\sigma (k)\frac{1}{2\pi }_1^1๐q\frac{2\alpha \sigma (q)}{\alpha ^2+(kq)^2}=\frac{1}{2\pi },$$
(14)
where $`\alpha `$ is an implicit function of $`\gamma :`$
$$\alpha =\gamma _1^1๐k\sigma (k).$$
(15)
In terms of the moments $`ฯต_m`$ defined by
$$ฯต_m=\left(\frac{\gamma }{\alpha }\right)^{m+1}_1^1๐kk^m\sigma (k),m=2,4.$$
(16)
the exact expression for $`g_3(\gamma ,\mathrm{})`$ is
$$\begin{array}{c}\frac{g_3(\gamma ,\mathrm{})}{n^3}=\frac{3}{2\gamma }ฯต_4^{}\frac{5ฯต_4}{\gamma ^2}\hfill \\ \hfill +\left(1+\frac{\gamma }{2}\right)ฯต_2^{}2\frac{ฯต_2}{\gamma }\frac{3ฯต_2ฯต_2^{}}{\gamma }+\frac{9ฯต_2^2}{\gamma ^2},\end{array}$$
(17)
where $`ฯต_m^{}`$ is the derivative of $`ฯต_m`$ with respect to $`\gamma .`$
Equations (14)โ(17) form a closed set of equations determining $`g_3(\gamma ,\mathrm{}).`$ We plot the result in Figs. 1 and 2.
For a numerical solution of the integral equation (14) we have used the Mathematica package NISolve NISolve . In Fig. 1 we compare the normalized function $`g_3(\gamma ,\mathrm{})/g_3(0,\mathrm{})`$ with that calculated for $`N=3`$, $`4`$, $`5`$, and $`6.`$ One can see that the convergence to the thermodynamic limit is quite slow. In Fig. 2 we compare $`g_3(\gamma ,\mathrm{})/n^3`$ with the experimental data from Ref. Tolra . The average value of $`\gamma ,`$ measured in Tolra , $`\gamma _\mathrm{m}0.45,`$ and the corresponding value of $`g_3(\gamma _\mathrm{m},\mathrm{})0.14`$ are shown as the dot inside the box which represents the experimental uncertainty according to Ref. Tolra : $`0.34<\gamma _\mathrm{m}<0.65`$ and $`0.05<g_3(\gamma _\mathrm{m},\mathrm{})<0.23.`$ The dashed lines show the asymptotic expressions given in Eqs. (12) and (13) for small and large $`\gamma ,`$ respectively. Evidently, the asymptotic expressions do not account for the observed value, and the exact expression (17) is needed.
Eq. (17) gives $`g_3(\gamma ,\mathrm{})/n^316\pi ^6/15\gamma ^6`$ when $`\gamma \mathrm{},`$ thus reproducing the asymptotic expression Eq. (13) taken at $`N=\mathrm{}.`$ To check this is a straightforward task since Eq. (14) admits a regular perturbative expansion with $`1/\alpha `$ being a small parameter. The opposite limit, $`\gamma 0,`$ of Eq. (17) is much more difficult to analyze since the kernel in Eq. (14) becomes singular when $`\alpha 0.`$ Beyond the leading order, the results were obtained only recently Wadati :
$$\sigma (k)\sqrt{1k^2}/2\pi \alpha +f(k),\alpha 0$$
(18)
where $`f(k)`$ is written explicitly in Ref. Wadati , see Eq. (4.11). For the moments Eq. (16) one gets $`ฯต_2\gamma (14\sqrt{\gamma }/3\pi )`$ and $`ฯต_42\gamma ^2(144\sqrt{\gamma }/15\pi ).`$ With these expressions for the moments, expansion (12) is reproduced.
A useful representation of our result (17) is the following approximate form
$`{\displaystyle \frac{16\pi ^1\gamma ^{\frac{1}{2}}+1.2656\gamma 0.2959\gamma ^{\frac{3}{2}}}{10.2262\gamma 0.1981\gamma ^{\frac{3}{2}}}},`$ $`0\gamma 1,`$ (19)
$`{\displaystyle \frac{0.7050.107\gamma +5.0810^3\gamma ^2}{1+3.41\gamma +0.903\gamma ^2+0.495\gamma ^3}},`$ $`1\gamma 7,`$ (20)
$`{\displaystyle \frac{16\pi ^6}{15\gamma ^6}}{\displaystyle \frac{9.435.40\gamma +\gamma ^2}{89.32+10.19\gamma +\gamma ^2}},`$ $`7\gamma 30.`$ (21)
The relative error of this approximation does not exceed $`210^3`$ for all values of $`\gamma `$ in the interval $`0\gamma 30`$, which we expect to be experimentally relevant. Note that the asymptotic result (13) for $`N=\mathrm{}`$ is nearly a factor of two larger than the exact result for $`\gamma =30`$.
The formal derivation of Eq. (17) is quite lengthy and will be given elsewhere CSZ . Here we shall only summarize the main steps: We are interested in the thermodynamic limit of the model, defined by Eq. (3). In this limit it is natural to use a field-theoretical approach rather than the first-quantized formalism we employed for the few-particle systems. The Lieb-Liniger model is completely integrable. In the field-theoretical language complete integrability implies the existence of an infinite family of conserved currents $`๐ฅ_n,`$ $`n=0,1,2,\mathrm{}`$. The integrals over space of these conserved currents commute with each other, including the Hamiltonian (4). There is thus an infinite set of mutually commuting operators $`J_n`$, $`n=0,1,2,\mathrm{},`$ for which the eigenfunctions and the spectrum are known exactly. A procedure for generating $`J_n`$ is discussed, for example, in Ref. KBI-93 . The first three currents generated by this method are $`๐ฅ_0=๐ฉ,`$ $`๐ฅ_1=๐ซ,`$ and $`๐ฅ_2=,`$ where $`๐ฉ,`$ $`๐ซ,`$ and $``$ are the number density of particles, the momentum density, and the Hamiltonian density, respectively. These three currents exist for non-integrable models as well, while the currents $`๐ฅ_n`$ for $`n>2`$ are specific to the LL model. To illustrate our approach we introduce the classical expression for the current $`๐ฅ_4`$ (its quantization will be discussed in the next paragraph):
$$\begin{array}{c}๐ฅ_4=\mathrm{\Psi }^{}_x^4\mathrm{\Psi }+c\mathrm{\Psi }^{}\mathrm{\Psi }^{}(_x\mathrm{\Psi })_x\mathrm{\Psi }+8c(_x\mathrm{\Psi }^{})\mathrm{\Psi }^{}\mathrm{\Psi }_x\mathrm{\Psi }\hfill \\ \hfill +c(_x\mathrm{\Psi }^{})(_x\mathrm{\Psi }^{})\mathrm{\Psi }\mathrm{\Psi }+2c^2(\mathrm{\Psi }^3)^{}\mathrm{\Psi }^3.\end{array}$$
(22)
The quantum-mechanical version of the last term in Eq. (22) contains $`\widehat{\mathrm{\Psi }}^{}(x)^3\widehat{\mathrm{\Psi }}(x)^3`$ in which we are interested. When averaged over the ground state, this term gives the desired correlation function $`g_3,`$ Eq. (1). Some of the remaining terms in (22) can be eliminated by using the Hellmann-Feynman theorem, which applies to all conserved quantities. This theorem was applied to the Hamiltonian of the LL model in Ref. Gangardt to calculate the correlation function $`\widehat{\mathrm{\Psi }}^{}(x)^2\widehat{\mathrm{\Psi }}(x)^2`$. In our case not all terms can be eliminated by use of the Hellmann-Feynman theorem. We have succeeded in obtaining additional identities CSZ by employing the long-distance asymptotics of the correlation functions of conserved currents known from the conformal limit of the theory.
The quantization of the classical current Eq. (22) is far from being straightforward. A naive quantum-mechanical version of the expression (22) contains terms involving derivatives of field operators $`\widehat{\mathrm{\Psi }}`$ and $`\widehat{\mathrm{\Psi }}^{}`$. Such terms are not well-defined in general. An explicit demonstration of this was considered in the paper Ol , where it was shown that the expectation value $`\widehat{\mathrm{\Psi }}^{}(0)_x^3\widehat{\mathrm{\Psi }}(x)`$ is discontinuous at $`x=0`$. The reason for this behavior can be seen from the structure of the Bethe-ansatz wave function (6), which has a cusp whenever the coordinates of any two particles coincide. This does not cause difficulties when one considers the momentum operator and the Hamiltonian itself, that is, when one works with the operators $`\widehat{\mathrm{\Psi }}^{}(x)_x\widehat{\mathrm{\Psi }}(x)`$ and $`\widehat{\mathrm{\Psi }}^{}(x)_x^2\widehat{\mathrm{\Psi }}(x).`$ However, for other operators we encounter difficulties which we have not been able to resolve within the continuum model (4). Therefore we have used an integrable lattice model, the so-called $`q`$-boson hopping model BBP-93 ; BIK-98 , which reduces to the Lieb-Liniger model (4) in the continuum limit. Within this model all higher (quantum) integrals of motion are defined unambiguously.
In conclusion, we have obtained an exact expression for the three-body correlation function Eq. (1) within the Lieb-Liniger model of a one-dimensional Bose gas. The explicit expression is given by Eq. (17) and Eqs. (19)-(21), and compared in Fig. 2 with a recent experiment. The finite-size effects are explored and a comparison with the thermodynamic limit is given in Fig. 1. We expect that our results will stimulate further experimental investigations of the local correlations in low-dimensional trapped atomic gases.
The authors would like to thank N.M. Bogoliubov for helpful discussions. M.B. Zvonarevโs work was supported by the Danish Technical Research Council via the Framework Programme on Superconductivity.
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# 1 Introduction
## 1 Introduction
One of the important tools for studying the properties of quark-gluon plasma (QGP) in ultrarelativistic heavy ion collisions is the analysis of a QCD jet production. The medium-induced energy loss of energetic partons, โjet quenchingโ, should be very different in the cold nuclear matter and QGP, resulting in many observable phenomena . Recent RHIC data on high-p<sub>T</sub> particle production (the suppression of hadron spectra and azimuthal back-to-back two-particle correlations, strong elliptic flow) are in agreement with the jet quenching hypothesis . At LHC, a new regime of heavy ion physics will be reached at $`\sqrt{s_{\mathrm{NN}}}=5.5A`$ TeV where hard and semi-hard particle production can stand out against the underlying soft events. The initial gluon densities in Pb+Pb reactions at LHC are expected to be significantly higher than those at RHIC, implying a stronger partonic energy loss, observable in various new channels .
In the most of available Monte-Carlo heavy ion event generators the medium-induced partonic rescattering and energy loss are either ignored or implemented insufficiently . Thus, in order to analyze RHIC data on high-p<sub>T</sub> hadron production, test the sensitivity of LHC observables to the QGP formation, and study the corresponding experimental capabilities of detectors, the development of adequate and fast Monte-Carlo models for jet quenching simulation is necessary.
In this paper we present the model of jet quenching and discuss its validation basing on RHIC data for high-$`p_T`$ hadron spectra. In Sect. 2 we give the physics frameworks of the model. Section 3 describes the event-by-event simulation procedure. In Sect. 4 the generalization of the model to the case of โfullโ heavy ion event (superposition of soft hydro-type state and hard multi-jets) is fulfilled. In Sect. 5 the efficiency of the model is demonstrated by means of the numerical analysis of hadron spectra at RHIC.
## 2 Physics frameworks of the model
The detailed description of physics frameworks of the developed model can be found in a number of our previous papers . Our approach bases on an accumulating energy loss, the gluon radiation being associated with each parton scattering in the expanding medium and includes the interference effect using the modified radiation spectrum $`dE/dl`$ as a function of decreasing temperature $`T`$. The basic kinetic integral equation for the energy loss $`\mathrm{\Delta }E`$ as a function of initial energy $`E`$ and path length $`L`$ has the form
$`\mathrm{\Delta }E(L,E)={\displaystyle \underset{0}{\overset{L}{}}}๐l{\displaystyle \frac{dP(l)}{dl}}\lambda (l){\displaystyle \frac{dE(l,E)}{dl}},{\displaystyle \frac{dP(l)}{dl}}={\displaystyle \frac{1}{\lambda (l)}}\mathrm{exp}\left(l/\lambda (l)\right),`$ (1)
where $`l`$ is the current transverse coordinate of a parton, $`dP/dl`$ is the scattering probability density, $`dE/dl`$ is the energy loss per unit length, $`\lambda =1/(\sigma \rho )`$ is in-medium mean free path, $`\rho T^3`$ is the medium density at the temperature $`T`$, $`\sigma `$ is the integral cross section for parton interaction in the medium.
The collisional energy loss due to elastic scattering with high-momentum transfer has been originally estimated by Bjorken in , and recalculated later in taking also into account the low-momentum transfer loss resulting mainly from the interactions with plasma collective modes. Since the latter process does not contribute much to the total collisional loss in comparison with high-momentum scattering (due to absence of large factor $`\mathrm{ln}(E/\mu _D)`$ where $`\mu _D`$ is the Debye screening mass) and, in numerical estimates it can be effectively โabsorbedโ by means of redefinition of minimum momentum transfer $`t_{\mathrm{min}}\mu _D^2`$ , we used the collisional part associated with high-momentum transfer only ,
$$\frac{dE}{dl}^{col}=\frac{1}{4T\lambda \sigma }\underset{\mu _D^2}{\overset{t_{\mathrm{max}}}{}}๐t\frac{d\sigma }{dt}t,$$
(2)
and the dominant contribution to the differential cross section
$$\frac{d\sigma }{dt}C\frac{2\pi \alpha _s^2(t)}{t^2}\frac{E^2}{E^2m_p^2},\alpha _s=\frac{12\pi }{(332N_f)\mathrm{ln}(t/\mathrm{\Lambda }_{QCD}^2)}$$
(3)
for scattering of a hard parton with energy $`E`$ and mass $`m_p`$ off the โthermalโ parton with energy (or effective mass) $`m_03TE`$. Here $`C=9/4,1,4/9`$ for $`gg`$, $`gq`$ and $`qq`$ scatterings respectively, $`\alpha _s`$ is the QCD running coupling constant for $`N_f`$ active quark flavors, and $`\mathrm{\Lambda }_{QCD}`$ is the QCD scale parameter which is of the order of the critical temperature, $`\mathrm{\Lambda }_{QCD}T_c200`$ MeV. The integrated cross section $`\sigma `$ is regularized by the Debye screening mass squared $`\mu _D^2(T)4\pi \alpha _sT^2(1+N_f/6)`$. The maximum momentum transfer $`t_{\mathrm{max}}=[s(m_p+m_0)^2][s(m_pm_0)^2]/s`$ where $`s=2m_0E+m_0^2+m_p^2`$.
There are several calculations of the inclusive energy distribution of medium-induced gluon radiation using Feyman multiple scattering diagrams. The relation between these approaches and their basic parameters has been discussed in detail in the recent writeup of the working group โJet Physicsโ for the CERN Yellow Report . We restrict ourselves to using BDMS formalism . In the BDMS frameworks, the strength of multiple scattering is characterized by the transport coefficient $`\widehat{q}=\mu _D^2/\lambda _g`$ ($`\lambda _g`$ is the gluon mean free path), which is related to the elastic scattering cross section $`\sigma `$ (3). In our simulations this strength is in fact regulated mainly by the initial QGP temperature $`T_0`$. Then the energy spectrum of coherent medium-induced gluon radiation and the corresponding dominant part of radiative energy loss of massless parton have the form :
$`{\displaystyle \frac{dE}{dl}}^{rad}={\displaystyle \frac{2\alpha _s(\mu _D^2)C_R}{\pi L}}{\displaystyle \underset{\omega _{\mathrm{min}}}{\overset{E}{}}}๐\omega \left[1y+{\displaystyle \frac{y^2}{2}}\right]\mathrm{ln}\left|\mathrm{cos}(\omega _1\tau _1)\right|,`$ (4)
$`\omega _1=\sqrt{i\left(1y+{\displaystyle \frac{C_R}{3}}y^2\right)\overline{\kappa }\mathrm{ln}{\displaystyle \frac{16}{\overline{\kappa }}}}\text{with}\overline{\kappa }={\displaystyle \frac{\mu _D^2\lambda _g}{\omega (1y)}},`$ (5)
where $`\tau _1=L/(2\lambda _g)`$, $`y=\omega /E`$ is the fraction of the hard parton energy carried away by the radiated gluon, and $`C_R=4/3`$ is the quark color factor. A similar expression for the gluon jet can be obtained by setting $`C_R=3`$ and proper by changing the factor in the square brackets in (4), see ref. . The integration (4) is carried out over all energies from $`\omega _{\mathrm{min}}=E_{\mathrm{LPM}}=\mu _D^2\lambda _g`$, the minimum radiated gluon energy in the coherent LPM regime, up to initial parton energy $`E`$. Note that we do not consider here possible effects of double parton scattering and thermal gluon absorption , which can be included in the model in the future.
The simplest generalization of the formula for a heavy quark of mass $`m_q`$ can be done by using the โdead-coneโ approximation :
$$\frac{dE}{dld\omega }|_{m_q0}=\frac{1}{(1+(\beta \omega )^{3/2})^2}\frac{dE}{dld\omega }|_{m_q=0},\beta =\left(\frac{\lambda }{\mu _D^2}\right)^{1/3}\left(\frac{m_q}{E}\right)^{4/3}.$$
(6)
One should mention the more recent developments on heavy quark energy loss calculations available in the literature , which can be also considered as further model improvements.
The medium is treated as a boost-invariant longitudinally expanding quark-gluon fluid, and partons as being produced on a hyper-surface of equal proper times $`\tau `$ . In order to simplify numerical calculations we omit here the transverse expansion and viscosity of the fluid using the well-known scaling solution obtained by Bjorken for a temperature and density of QGP at $`T>T_c200`$ MeV:
$$\epsilon (\tau )\tau ^{4/3}=\epsilon _0\tau _0^{4/3},T(\tau )\tau ^{1/3}=T_0\tau _0^{1/3},\rho (\tau )\tau =\rho _0\tau _0.$$
(7)
The internal model parameters are the initial conditions for the QGP formation expected for central Au+Au (Pb+Pb) collisions at RHIC (LHC): $`\tau _0`$, $`T_0`$ and $`N_f`$. For non-central collisions and for other beam atomic numbers we suggest the proportionality of the initial energy density $`\epsilon _0`$ to the ratio of nuclear overlap function and effective transverse area of nuclear overlapping .
Note that using other scenarios of QGP space-time evolution for the Monte-Carlo implementation of the model is also envisaged. In fact, the influence of the transverse flow, as well as that of the mixed phase at $`T=T_c`$, on the intensity of jet rescattering (which is a strongly increasing function of $`T`$) has been found to be inessential for high initial temperatures $`T_0T_c`$. On the contrary, the presence of QGP viscosity slows down the cooling rate, that implies a jet parton spending more time in the hottest regions of the medium. As a result the rescattering intensity increases, i.e., in fact an effective temperature of the medium gets higher as compared with the perfect QGP case. We also do not take into account here the probability of jet rescattering in nuclear matter, because the intensity of this process and corresponding contribution to the total energy loss are not significant due to much smaller energy density in a โcoldโ nucleus.
Another important element of the model is the angular spectrum of in-medium gluon radiation. Since the detailed calculation of the angular spectrum of emitted gluons is rather sophisticated and model-dependent , the simple parameterization of gluon angular distribution over the emission angle $`\theta `$ was used:
$$\frac{dN^g}{d\theta }\mathrm{sin}\theta \mathrm{exp}\left(\frac{(\theta \theta _0)^2}{2\theta _0^2}\right),$$
(8)
where $`\theta _05^0`$ is the typical angle of the coherent gluon radiation as estimated in . Other parameterizations are also envisaged.
## 3 Simulation procedure
The model has been constructed as the fast Monte-Carlo event generator PYQUEN (PYthia QUENched), and the corresponding Fortran routine PYQUEN is available via Internet . The routine is implemented as a modification of the standard PYTHIA$`\mathrm{\_}`$6.2.\* jet event .
The following event-by-event Monte-Carlo simulation procedure is applied.
$``$ Generation of the initial parton spectra with PYTHIA (fragmentation off).
$``$ Generation of the jet production vertex at the impact parameter $`b`$ according to the distribution
$$\frac{dN^{\mathrm{jet}}}{d\psi dr}(b)=\frac{T_A(r_1)T_A(r_2)}{T_{AA}(b)},T_{AA}(b)=\underset{0}{\overset{2\pi }{}}๐\psi \underset{0}{\overset{r_{max}}{}}r๐rT_A(r_1)T_A(r_2),$$
(9)
where $`r_{1,2}(b,r,\psi )`$ are the distances between the nucleus centers and the jet production vertex $`V(r\mathrm{cos}\psi ,r\mathrm{sin}\psi )`$; $`r_{max}(b,\psi )R_A`$ is the maximum possible transverse distance $`r`$ from the nuclear collision axis to $`V`$; $`R_A`$ is the radius of the nucleus $`A`$; $`T_A(๐ซ)=A\rho _A(๐ซ,z)๐z`$ is the nuclear thickness function with nucleon density distribution $`\rho _A(๐ซ,z)`$; $`T_{AA}(b)`$ is the nuclear overlap function (see ref. for detailed nuclear geometry explanations).
$``$ Calculation of scattering cross section $`\sigma =๐t๐\sigma /๐t`$ (3).
$``$ Generation of the displacement between $`i`$-th and $`(i+1)`$-th scatterings, $`l_i=(\tau _{i+1}\tau _i)`$:
$$\frac{dP}{dl_i}=\lambda ^1(\tau _{i+1})\mathrm{exp}(\underset{0}{\overset{l_i}{}}\lambda ^1(\tau _i+s)๐s),\lambda ^1(\tau )=\sigma (\tau )\rho (\tau ),$$
(10)
and calculation of the corresponding transverse distance, $`l_ip_T/E`$.
$``$ Reducing the parton energy by collisional and radiative loss per each $`i`$-th scattering:
$$\mathrm{\Delta }E_{\mathrm{tot},i}=\mathrm{\Delta }E_{\mathrm{col},i}+\mathrm{\Delta }E_{\mathrm{rad},i},$$
(11)
where the collisional part is calculated in the high-momentum transfer approximation (3),
$$\mathrm{\Delta }E_{\mathrm{col},i}=\frac{t_i}{2m_0},$$
(12)
and the energy of a radiated gluon, $`\omega _i=\mathrm{\Delta }E_{\mathrm{rad},i}`$, is generated according to (4) and (6) :
$`{\displaystyle \frac{dI}{d\omega }}|_{m_q=0}={\displaystyle \frac{2\alpha _s(\mu _D^2)\lambda C_R}{\pi L\omega }}\left[1y+{\displaystyle \frac{y^2}{2}}\right]\mathrm{ln}\left|\mathrm{cos}(\omega _1\tau _1)\right|,{\displaystyle \frac{dI}{d\omega }}|_{m_q0}={\displaystyle \frac{1}{(1+(\beta \omega )^{3/2})^2}}{\displaystyle \frac{dI}{d\omega }}|_{m_q=0}.`$ (13)
$``$ Calculation of the parton transverse momentum kick due to elastic scattering $`i`$:
$$\mathrm{\Delta }k_{t,i}^2=(E\frac{t_i}{2m_{0i}})^2(p\frac{E}{p}\frac{t_i}{2m_{0i}}\frac{t_i}{2p})^2m_p^2.$$
(14)
$``$ Formation of the additional (in-medium emitted) gluon with the energy $`\omega _i`$ and the emission angle $`\theta _i`$ relative to the parent parton determined according to the parameterization (8).
$``$ Halting the rescattering if (1) the parton escapes the dense zone, or (2) QGP cools down to $`T_c=200`$ MeV, or (3) the parton loses so much energy that its $`p_T(\tau )`$ drops below $`2T(\tau )`$.
$``$ At the end of each event, adding new (in-medium emitted) gluons to the PYTHIA parton list and rearrangement of partons to update string formation.
$``$ Formation of the final state particles by PYTHIA (fragmentation on).
## 4 Extension of the model to simulate full heavy ion event
The full heavy ion event is simulated as a superposition of soft hydro-type state and hard multi-jets. The simple approximation of hadronic liquid at โfreeze-outโ stage has been used to treat soft part of the event giving the final hadron spectrum as a superposition of a thermal distribution and a collective flow .
1. The 4-momentum $`p_\mu ^{}`$ of a hadron of mass $`m`$ was generated at random in the rest frame of a liquid element in accordance with the isotropic Boltzmann distribution
$`f(E^{})E^{}\sqrt{E^2m^2}\mathrm{exp}(E^{}/T_f),1<\mathrm{cos}\theta ^{}<1,0<\varphi ^{}<2\pi ,`$ (15)
where $`E^{}=\sqrt{p^2+m^2}`$ is the energy of the hadron, and the polar angle $`\theta ^{}`$ and the azimuthal angle $`\varphi ^{}`$ specify the direction of its motion in the rest frame of the liquid element.
2. The spatial position of a liquid element and its local 4-velocity $`u_\mu `$ were generated at random in accordance with phase space and the character of motion of the fluid:
$`f(r)=2r/R_f^2(0<r<R_f),f(\eta )\mathrm{exp}\left[(\eta Y_L^{\mathrm{max}})^2/2(Y_L^{\mathrm{max}})^2\right],0<\mathrm{\Phi }<2\pi ,`$
$`u_r=\mathrm{sinh}Y_T^{\mathrm{max}}{\displaystyle \frac{r}{\sqrt{R_f(b)R_f(b=0)}}},u_t=\sqrt{1+u_r^2}\mathrm{cosh}\eta ,u_z=\sqrt{1+u_r^2}\mathrm{sinh}\eta ,`$ (16)
where $`R_f`$ is the final transverse radius of the system in a given direction. Freeze-out parameters of the model are kinetic freeze-out temperature $`T_f`$ and maximum longitudinal, $`Y_L^{\mathrm{max}}`$, and transverse, $`Y_T^{\mathrm{max}}`$, collective flow rapidities.
3. Further, boost of the hadron 4-momentum in the c.m. frame of the event was calculated:
$`p_x`$ $`=`$ $`p^{}\mathrm{sin}\theta ^{}\mathrm{cos}\varphi ^{}+u_r\mathrm{cos}\mathrm{\Phi }\left[E^{}+{\displaystyle \frac{(u^ip^i)}{u_t+1}}\right]`$
$`p_y`$ $`=`$ $`p^{}\mathrm{sin}\theta ^{}\mathrm{sin}\varphi ^{}+u_r\mathrm{sin}\mathrm{\Phi }\left[E^{}+{\displaystyle \frac{(u^ip^i)}{u_t+1}}\right]`$
$`p_z`$ $`=`$ $`p^{}\mathrm{cos}\theta ^{}+u_z\left[E^{}+{\displaystyle \frac{(u^ip^i)}{u_t+1}}\right]`$
$`E`$ $`=`$ $`E^{}u_t+(u^ip^i),`$ (17)
where
$`(u^ip^i)`$ $`=`$ $`u_rp^{}\mathrm{sin}\theta ^{}\mathrm{cos}(\mathrm{\Phi }\varphi ^{})+u_zp^{}\mathrm{cos}\theta ^{}.`$ (18)
Anisotropic flow is introduced here under simple assumption that the spatial ellipticity of โfreeze-outโ region, $`ฯต=y^2x^2/y^2+x^2`$, is directly related to the ellipticity of the system formed in the region of the initial overlap of nuclei, $`ฯต_0=b/2R_A`$. This โscalingโ enables one to avoid introducing additional parameters and, at the same time, leads to an azimuthal anisotropy of generated particles due to dependence of transverse radius $`R_f(b)`$ on the angle $`\mathrm{\Phi }`$ :
$$R_f(b)=R_f(b=0)\mathrm{min}\{\sqrt{1ฯต_0^2\mathrm{sin}^2\mathrm{\Phi }}+ฯต_0\mathrm{cos}\mathrm{\Phi },\sqrt{1ฯต_0^2\mathrm{sin}^2\mathrm{\Phi }}ฯต_0\mathrm{cos}\mathrm{\Phi }\}.$$
(19)
Obtained in such a way azimuthal distribution of particles is described well by the elliptic form for the domain of reasonable impact parameter values.
The mean total particle multiplicity in central Au+Au (Pb+Pb) collisions at RHIC (LHC) is the input parameter of the model (instead of $`R_f(b=0)`$ we put $`R_A`$ here for simplicity), the total multiplicity for other centralities and atomic numbers being assumed to be proportional to the number of nucleons-participants. We also set the Poisson multiplicity distribution and the following particle ratios:
$$\pi ^\pm :K^\pm :p^\pm =24:6:1,\pi ^\pm :\pi ^0=2:1,K^\pm :K^0=1:1,p:n=1:1.$$
The hard part of the event includes PYTHIA/PYQUEN hadronic jets generated according to the binomial distribution. The mean number of jets produced in AA events at a given $`b`$ is proportional to the number of binary nucleon-nucleon sub-collisions and determined as
$$\overline{N_{AA}^{\mathrm{jet}}}(b,\sqrt{s})=T_{AA}(b)\underset{p_T^{\mathrm{min}}}{}๐p_T^2๐y\frac{d\sigma _{pp}^{\mathrm{hard}}(p_T,\sqrt{s})}{dp_T^2dy},$$
(20)
where $`d\sigma _{pp}^{\mathrm{hard}}(p_T,\sqrt{s})/dp_T^2dy`$ is the cross section of corresponding hard process in $`pp`$ collisions (at the same c.m.s. energy, $`\sqrt{s}`$, of colliding beams) with the minimum transverse momentum transfer $`p_T^{\mathrm{min}}`$. The latter is another input parameter of the model. In the frameworks of our approximation, partons produced in (semi)hard processes with the momentum transfer less than $`p_T^{\mathrm{min}}`$ are considered as being โthermalizedโ, so their hadronization products are included in a soft part of the event โautomaticallyโ.
Note that we can expect some adequate results only for central and semi-central collisions, but not for very peripheral collisions ($`b2R_A`$) where the hydro-type description is not applicable. Besides, the very forward rapidity region (where other dynamical effects can be important) is beyond our treatment here.
The extended in such a way jet quenching model has been constructed as the fast Monte-Carlo event generator, and the corresponding Fortran code is also available via Internet .
Let us remind in the end of this section, that ideologically our approximation is similar to the model of Hirano and Nara . The difference is that we concentrate here on the detailed simulation of the parton multiple scattering in a QCD-medium (the scattering-by-scattering generation of parton path length and energy loss in an expanding QGP, taking into account the collisional loss, Lund string fragmentation model both for hard partons and in-medium emitted gluons, etc.), while the treatment of the hydrodynamic part in is much more detailed than in our simple (and therefore fast) simulation procedure.
## 5 Validation of the model at RHIC, $`\sqrt{s}=200A`$ GeV
In order to demonstrate the efficiency of the model, the jet quenching pattern in Au+Au collisions at RHIC was considered. The comparison of calculated and experimentally measured pseudorapidity $`\eta `$ and transverse momentum $`p_T`$ spectra of hadrons together with their dependence of event centrality allows the optimization of the model and specification of main model parameters.
The PHOBOS data on $`\eta `$-spectra of charged hadrons have been analyzed to fix the particle density in the mid-rapidity region and the maximum longitudinal flow rapidity, $`Y_L^{\mathrm{max}}=3.5`$. For the calculation of (multi)jet production cross section, we used the factor $`K=2`$ taking into account higher order corrections of perturbative QCD. The rest of the model parameters have been obtained by fitting PHENIX data on $`p_T`$-spectra of neutral pions : the kinetic freeze-out temperature $`T_f=100`$ MeV, maximum transverse flow rapidity $`Y_T^{\mathrm{max}}=1.25`$ and minimum transverse momentum transfer of โnon-thermalizedโ hard process $`p_T^{\mathrm{min}}=2.8`$ GeV/$`c`$. It was found that the nuclear modification of the hardest domain of $`p_T`$-spectrum ($`\begin{array}{c}>\hfill \\ \hfill \end{array}5`$ GeV/$`c`$) is determined in our case only by the intensity of the medium-induced parton rescattering. This fact allows us to extract from the data initial conditions of the QGP formation independently on other input parameters: the initial temperature $`T_0=500`$ MeV, the formation time $`\tau _0=0.4`$ fm/$`c`$ and the number of active quark flavours $`N_f=2`$. We will see below, that setting model parameters as it was described above, makes it possible to reproduce the main features of jet quenching pattern at RHIC: the $`p_T`$โdependence of the nuclear modification factor $`R_{AA}`$ and two-particle azimuthal correlation function $`C(\mathrm{\Delta }\phi )`$.
Figure 1 shows $`\eta `$-distribution of charged hadrons in Au+Au collisions for different centrality sets. The good fit of PHOBOS data is achieved excepting very forward rapidities. The $`p_T`$-distributions of $`\pi ^0`$-mesons obtained at PHENIX is also well reproduced by our calculations, even for relatively peripheral collisions (Figure 2).
Figure 3 shows the nuclear modification factor $`R_{AA}`$ for neutral pions, which is defined as a ratio of particle yields in $`AA`$ and $`pp`$ collisions normalized on the number of binary nucleon-nucleon sub-collisions:
$$R_{AA}=\frac{d\sigma _{AA}^{\pi ^0}/dp_T}{T_{AA}(b)\sigma _{\mathrm{in}}d\sigma _{pp}^{\pi ^0}/dp_T},$$
(21)
where $`\sigma _{\mathrm{in}}=42`$ mb is the inelastic non-diffractive $`pp`$ cross section at $`\sqrt{s}=200`$ GeV. In the absence of medium-induced effects in the mid-rapidity region it should be $`R_{AA}=1`$ for high enough $`p_T(\begin{array}{c}>\hfill \\ \hfill \end{array}2`$ GeV/$`c`$). Such value of $`R_{AA}sim1`$ has been observed so for d+Au and peripheral Au+Au collisions, but not for central and semi-central Au+Au events, where $`R_{AA}<1`$ up to maximum measured transverse momenta $`p_T10`$ GeV/$`c`$. One can see from Figure 3 that our model calculations reproduce $`p_T`$โ and centrality dependences of $`R_{AA}`$ quite well.
Another important tool to verify jet quenching is two-particle azimuthal correlation function $`C(\mathrm{\Delta }\phi )`$ โ the distribution over an azimuthal angle of high-$`p_T`$ hadrons in the event with $`2`$ GeV/$`c<p_T<p_T^{\mathrm{trig}}`$ relative to that for the hardest โtriggerโ particle with $`p_T^{\mathrm{trig}}>4`$ GeV/$`c`$. Figure 4 presents $`C(\mathrm{\Delta }\phi )`$ in $`pp`$ and in central Au+Au collisions (data from STAR ). Clear peaks in $`pp`$ collisions at $`\mathrm{\Delta }\phi =0`$ and $`\mathrm{\Delta }\phi =\pi `$ indicate a typical dijet event topology. Note that almost the same pattern has been observed in d+Au and peripheral Au+Au collisions. However, for most central Au+Au collisions the peak near $`\pi `$ disappears. It can be interpreted as the observation of monojet events due to the โabsorptionโ of one of the jets in a dense medium. Such event configuration corresponds to the situation when the dijet production vertex is close to the surface of the nuclear overlap region: then one partonic jet can escape the medium almost without re-interactions and then go to detectors, while second jet loses most of its initial energy due to a large number of rescatterings and therefore becomes unobservable . Figure 4 demonstrates that measured suppression of azimuthal back-to-back correlations is well reproduced by our model (the same procedure of uncorrelated background subtraction as in was applied).
We leave beyond the scope of this paper the analysis of such important RHIC observables as the azimuthal anisotropy and particle ratios. Since these observables are very sensitive to the soft physics, in order to study them a more careful treatment of low-$`p_T`$ particle production than our simple approach is needed (the detailed description of space-time structure of freeze-out region, resonance decays, etc.) For example, our model can reproduce experimentally measured $`p_T`$-dependence of the coefficient of azimuthal anisotropy $`v_2`$ (the second harmonic of Fourier decomposition of particle azimuthal distribution) qualitatively, giving rapid hydrodynamical growth up to $`p_T3`$ GeV/$`c`$ with the subsequent saturation. However, the model calculations significantly underestimate the data at $`p_T<2`$ GeV/$`c`$. A solution of baryon-to-meson ratio โpuzzleโ is also beyond our consideration here. Further development of our model with special emphasis on the more detailed description of low-$`p_T`$ particle production is planed for the future.
## 6 Conclusions
The model of jet quenching in ultrarelativistic heavy ion collisions has been developed. It includes the generation of the hard parton production vertex according to the realistic nuclear geometry, rescattering-by-rescattering simulation of the parton path length in a dense matter, radiative and collisional energy loss per rescattering, final hadronization with the Lund string fragmentation model for hard partons and in-medium emitted gluons. The model is the fast Monte-Carlo tool implemented to modify a standard PYTHIA jet event. The model has been generalized to the case of the โfullโ heavy ion event (the superposition of soft, hydro-type state and hard multi-jets) using a simple and fast simulation procedure for soft particle production.
The efficiency of the model is demonstrated basing on the numerical analysis of high-$`p_T`$ hadron production in Au+Au collisions at RHIC. The good fit of experimental data on $`\eta `$โ and $`p_T`$โ spectra of hadrons for different event centralities is achieved. The model is capable of reproducing main features of the jet quenching pattern at RHIC: the $`p_T`$ dependence of the nuclear modification factor $`R_{AA}`$, and the suppression of azimuthal back-to-back correlations. The further development of the model focusing on a more detailed description of low-$`p_T`$ particle production is planed for the future.
Acknowledgments. Discussions with A.I. Demianov, Yu.L. Dokshitzer, A. Morsch, S.V. Petrushanko, C. Roland, L.I. Sarycheva, J. Schukraft, C.Yu. Teplov, I.N. Vardanyan, I. Vitev, B. Wyslouch, B.G. Zakharov and G.M. Zinovjev and are gratefully acknowledged. This work is supported by grant N 04-02-16333 of Russian Foundation for Basic Research.
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# Coupled nonlinear Schrรถdinger systems with potentials
## 1. Introduction
Very recently, different authors focused their attention on coupled nonlinear Schrรถdinger systems which describe physical phenomena such as the propagation in birefringent optical fibers, Kerr-like photorefractive media in optics and Bose-Einstein condensates.
First of all, let us recall that, in the last twenty years, motivated by the study of the propagation of pulse in nonlinear optical fiber, the nonlinear Schrรถdinger equation,
$$\mathrm{\Delta }u+u=u^3\mathrm{in}^3,$$
has been faced by many authors. It has been proved the existence of the least energy solution (ground state solution), which is radial with respect to some point, positive and exponentially decaying with its first derivatives at infinity. Moreover there are also many papers about the semiclassical states for the nonlinear Schrรถdinger equation with the presence of potentials
$$\epsilon ^2\mathrm{\Delta }u+V(x)u=u^3\mathrm{in}^3,$$
giving sufficient and necessary conditions to the existence of solutions concentrating in some points, and recently, in set with non-zero dimension, (see e.g. ).
However, by I.P. Kaminow , we know that single-mode optical fibers are not really โsingle-modeโ, but actually bimodal due to the presence of birefringence. This birefringence can deeply influence the way in which an optical evolves during the propagation along the fiber. Indeed, it can occur that the linear birefringence makes a pulse split in two, while nonlinear birefringent traps them together against splitting. C.R. Menyuk showed that the evolution of two orthogonal pulse envelopes in birefringent optical fibers is governed by the following coupled nonlinear Schrรถdinger system
(1.1)
$$\{\begin{array}{cc}i\varphi _t+\varphi _{xx}+|\varphi |^2\varphi +\beta |\psi |^2\varphi =0,\hfill & \\ i\psi _t+\psi _{xx}+|\psi |^2\psi +\beta |\varphi |^2\psi =0,\hfill & \end{array}$$
with $`\beta `$ positive constant depending on the anisotropy of the fibers. System (1.1) is also important for industrial applications in fiber communications systems and all-optical switching devices . If one seeks for standing wave solutions of (1.1), namely solutions of the form
$$\varphi (x,t)=e^{iw_1^2t}u(x)\mathrm{and}\psi (x,t)=e^{iw_2^2t}v(x),$$
then (1.1) becomes
(1.2)
$$\{\begin{array}{cc}u_{xx}+u=|u|^2u+\beta |v|^2u\hfill & \mathrm{in},\hfill \\ v_{xx}+w^2v=|v|^2v+\beta |u|^2v\hfill & \mathrm{in},\hfill \end{array}$$
with $`w^2=w_2^2/w_1^2`$. Finally we want to recall that (1.2) describes also other physical phenomena, such as Kerr-like photorefractive media in optics, (cf. ).
Problem (1.2), in a more general situation and also in higher dimension, has been studied by R. Cipolatti & W. Zumpichiatti . By concentration compactness arguments, they prove the existence and the regularity of the ground state solutions $`(u,v)(0,0)`$. Later on, in two very recent papers, T.C. Lin & J. Wei and L.A. Maia, E. Montefusco & B. Pellacci deal with problem (1.2), also in the multidimensional case, and, among other results, they prove the existence of least energy solutions of the type $`(u,v)`$, with $`u,v>0`$. Moreover T.C. Lin & J. Wei prove that, if $`\beta <0`$, then the ground state solution for (1.2) does not exist. We refer to all these papers and to references therein for more complete informations about (1.2).
Another motivation to the study of coupled Schrรถdinger systems arises from the Hartree-Fock theory for the double condensate, that is a binary mixture of Bose-Einstein condensates in two different hyperfine states $`|1`$ and $`|2`$ (cf. ). Indeed these phenomena are governed by the following system:
(1.3)
$$\{\begin{array}{cc}\epsilon ^2\mathrm{\Delta }u+\lambda _1u=\mu _1u^3+\beta uv^2\hfill & \mathrm{in}\mathrm{\Omega },\hfill \\ \epsilon ^2\mathrm{\Delta }v+\lambda _2v=\mu _2v^3+\beta u^2v\hfill & \mathrm{in}\mathrm{\Omega },\hfill \\ u,v>0\hfill & \mathrm{in}\mathrm{\Omega },\hfill \\ u=v=0\hfill & \mathrm{on}\mathrm{\Omega },\hfill \end{array}$$
where $`\mathrm{\Omega }`$ is a bounded domain of $`^3`$. Physically, $`u`$ and $`v`$ represent the corresponding condensate amplitudes, $`\epsilon ^2=\frac{\mathrm{}^2}{2m}`$, with $`\mathrm{}`$ the Planck constant and $`m`$ the atom mass. Moreover $`\mu _j=(N_j1)U_{jj}`$ and $`\beta =N_2U_{12}`$, with $`N_j1`$ a fixed number of atoms in the hyperfine state $`|j`$, and $`U_{ij}=4\pi \frac{\mathrm{}^2}{m}a_{ij}`$, where $`a_{jj}`$โs and $`a_{12}`$ are the intraspecies and interspecies scattering lengths. Besides, by E. Timmermans , we infer that $`\mu _j=\mu _j(x)`$ represents a chemical potential. For more informations about (1.3), see and references therein.
T.C. Lin & J. Wei, in , studied problem (1.3) with $`\lambda _1,\lambda _2,\mu _1,\mu _2`$ positive constant and they proved that if $`\beta <\sqrt{\mu _1\mu _2}`$, for $`\epsilon `$ sufficiently small, (1.3) has a least energy solution $`(u_\epsilon ,v_\epsilon )`$. Moreover, they distinguished two cases: the attractive case and the repulsive one. In the attractive case, which occurs whenever $`\beta >0`$, $`u_\epsilon `$ and $`v_\epsilon `$ concentrate respectively in $`Q_\epsilon `$ and $`Q_\epsilon ^{}`$, with
$$\frac{|Q_\epsilon Q_\epsilon ^{}|}{\epsilon }0,\mathrm{as}\epsilon 0.$$
Precisely they proved that
$`d(Q_\epsilon ,\mathrm{\Omega })\underset{Q\mathrm{\Omega }}{\mathrm{max}}d(Q,\mathrm{\Omega }),`$
$`d(Q_\epsilon ^{},\mathrm{\Omega })\underset{Q\mathrm{\Omega }}{\mathrm{max}}d(Q,\mathrm{\Omega }).`$
In the repulsive case, that is when $`\beta <0`$, the concentration points $`Q_\epsilon `$ and $`Q_\epsilon ^{}`$ satisfy the following condition:
$$\phi (Q_\epsilon ,Q_\epsilon ^{})\underset{(Q,Q^{})\mathrm{\Omega }^2}{\mathrm{max}}\phi (Q,Q^{}),$$
where
$$\phi (Q,Q^{})=\mathrm{min}\{\sqrt{\lambda _1}|QQ^{}|,\sqrt{\lambda _2}|QQ^{}|,\sqrt{\lambda _1}d(Q,\mathrm{\Omega }),\sqrt{\lambda _2}d(Q^{},\mathrm{\Omega })\}.$$
In particular
$$\frac{|Q_\epsilon Q_\epsilon ^{}|}{\epsilon }\mathrm{},\mathrm{as}\epsilon 0.$$
Motivated by these results and by the fact that we know that $`\mu _j`$ may be not constants (cf. ), in this paper we study the following problem:
($`๐ซ_\epsilon `$)
$$\{\begin{array}{cc}\epsilon ^2\mathrm{\Delta }u+J_1(x)u=J_2(x)u^3+\beta uv^2\hfill & \mathrm{in}\mathrm{\Omega },\hfill \\ \epsilon ^2\mathrm{\Delta }v+K_1(x)v=K_2(x)v^3+\beta u^2v\hfill & \mathrm{in}\mathrm{\Omega },\hfill \\ u,v>0\hfill & \mathrm{in}\mathrm{\Omega },\hfill \\ u=v=0\hfill & \mathrm{on}\mathrm{\Omega },\hfill \end{array}$$
with $`\mathrm{\Omega }^3`$, possibly unbounded and with smooth boundary, and with $`\beta <0`$, namely in the repulsive case. We will show that the presence of the potentials change drastically the situation with respect to the case with positive constants for what concerns the location of peaks, but, in some sense, not the repulsive nature of the problem. In fact, with suitable assumptions on the potentials, for $`\epsilon `$ sufficiently small, we will find solutions $`(u_\epsilon ,v_\epsilon )`$ of ($`๐ซ_\epsilon `$), even if not of least energy, concentrating respectively on $`Q_\epsilon `$ and $`Q_\epsilon ^{}`$ which tend toward the same point, determined by the potentials, as $`\epsilon 0`$, but with the property that the distance between them divided by $`\epsilon `$ diverges (see Remark 1.2)
Up to our knowledge, in this paper we give a first existence result of concentrating solutions for problem ($`๐ซ_\epsilon `$), in presence of potentials.
On the potentials $`J_i`$ and $`K_i`$ we will do the following assumptions:
* for $`i=1,2`$, $`J_iC^1(\mathrm{\Omega },)`$, $`J_i`$ and $`DJ_i`$ are bounded; moreover,
$$J_i(x)C>0\text{for all }x\mathrm{\Omega };$$
* for $`i=1,2`$, $`K_iC^1(\mathrm{\Omega },)`$, $`K_i`$ and $`DK_i`$ are bounded; moreover,
$$K_i(x)C>0\text{for all }x\mathrm{\Omega }.$$
Without lost of generality, we can suppose that there exists $`\epsilon _0>0`$, such that $`\mathrm{\Omega }_0:=\mathrm{\Omega }(\mathrm{\Omega }\epsilon _0e_1)\mathrm{}`$, where $`e_1=(1,0,0)`$.
Let us introduce an auxiliary function which will play a crucial role in the study of ($`๐ซ_\epsilon `$). Let $`\mathrm{\Gamma }:\mathrm{\Omega }_0`$ be a function so defined:
(1.4)
$$\mathrm{\Gamma }(Q)=J_1(Q)^{\frac{1}{2}}J_2(Q)^1+K_1(Q)^{\frac{1}{2}}K_2(Q)^1.$$
Let us observe that by (J) and (K), $`\mathrm{\Gamma }`$ is well defined.
Our main result is:
###### Theorem 1.1.
Suppose (J) and (K) and $`\beta <0`$. Let $`Q_0\mathrm{\Omega }_0`$ be an isolated local strict minimum or maximum of $`\mathrm{\Gamma }`$. There exists $`\overline{\epsilon }>0`$ such that if $`0<\epsilon <\overline{\epsilon }`$, then ($`๐ซ_\epsilon `$) possesses a solution $`(u_\epsilon ,v_\epsilon )`$ such that $`u_\epsilon `$ concentrates in $`Q_\epsilon `$ with $`Q_\epsilon Q_0`$, as $`\epsilon 0`$, and $`v_\epsilon `$ concentrates in $`Q_\epsilon ^{}`$ with $`Q_\epsilon ^{}Q_0`$, as $`\epsilon 0`$.
###### Remark 1.2.
Let us observe that, by the proof, it will be clear that, even if
$$|Q_\epsilon Q_\epsilon ^{}|0,\mathrm{as}\epsilon 0,$$
we have
$$\frac{|Q_\epsilon Q_\epsilon ^{}|}{\epsilon }\mathrm{},\mathrm{as}\epsilon 0.$$
Let us present how Theorem 1.1 becomes in some particular situations.
Let $`H:\mathrm{\Omega }`$ satisfying the assumption:
* $`HC^1(\mathrm{\Omega },)`$, $`H`$ and $`DH`$ are bounded; moreover,
$$H(x)C>0\text{for all }x\mathrm{\Omega }.$$
###### Corollary 1.3.
Suppose (H) and $`\beta <0`$. Suppose, moreover, that we are in one of the following situations:
* all the potentials $`J_i`$ and $`K_i`$ coincide with $`H`$;
* there exists $`i_0=1,2`$ such that $`J_{i_0}H`$ and $`K_{i_0}H`$, for $`i=i_0`$, while $`J_i`$ and $`K_i`$ are constant for $`ii_0`$;
* all the potentials $`J_i`$ and $`K_i`$ are constant, except only one, which coincides with $`H`$.
Let $`Q_0\mathrm{\Omega }_0`$ be an isolated local strict minimum or maximum of $`H`$. There exists $`\overline{\epsilon }>0`$ such that if $`0<\epsilon <\overline{\epsilon }`$, then ($`๐ซ_\epsilon `$) possesses a solution $`(u_\epsilon ,v_\epsilon )`$ such that $`u_\epsilon `$ concentrates in $`Q_\epsilon `$ with $`Q_\epsilon Q_0`$, as $`\epsilon 0`$, and $`v_\epsilon `$ concentrates in $`Q_\epsilon ^{}`$ with $`Q_\epsilon ^{}Q_0`$, as $`\epsilon 0`$.
###### Remark 1.4.
If, instead of $`\beta `$ constant, we consider $`\beta C^1(\mathrm{\Omega },)`$, bounded and bounded above by a negative constant, then we have exactly the same results.
Finally, we want to observe that we can treat also a more general problem than ($`๐ซ_\epsilon `$). Let us consider, indeed,
($`\overline{๐ซ}_\epsilon `$)
$$\{\begin{array}{cc}\epsilon ^2\mathrm{\Delta }u+J_1(x)u=J_2(x)u^{2p1}+\beta u^{p1}v^p\hfill & \mathrm{in}\mathrm{\Omega },\hfill \\ \epsilon ^2\mathrm{\Delta }v+K_1(x)v=K_2(x)v^{2p1}+\beta u^pv^{p1}\hfill & \mathrm{in}\mathrm{\Omega },\hfill \\ u,v>0\hfill & \mathrm{in}\mathrm{\Omega },\hfill \\ u=v=0\hfill & \mathrm{on}\mathrm{\Omega },\hfill \end{array}$$
with $`\mathrm{\Omega }^N`$, possibly unbounded and with smooth boundary, $`N3`$, $`2<2p<2N/(N2)`$ and with $`\beta <0`$.
Also in this case, without lost of generality, we can suppose that there exists $`\epsilon _0>0`$, such that $`\overline{\mathrm{\Omega }}_0:=\mathrm{\Omega }(\mathrm{\Omega }\epsilon _0\overline{e}_1)\mathrm{}`$, where $`\overline{e}_1=(1,0,\mathrm{},0)^N`$.
Let us define now $`\overline{\mathrm{\Gamma }}:\overline{\mathrm{\Omega }}_0`$ be a function so defined:
$$\overline{\mathrm{\Gamma }}(Q)=J_1(Q)^{\frac{p}{p1}\frac{N}{2}}J_2(Q)^{\frac{1}{p1}}+K_1(Q)^{\frac{p}{p1}\frac{N}{2}}K_2(Q)^{\frac{1}{p1}}.$$
In this case, Theorem 1.1 becomes:
###### Theorem 1.5.
Let $`N3`$ and $`2<2p<2N/(N2)`$. Suppose (J) and (K) and $`\beta <0`$. Let $`Q_0\overline{\mathrm{\Omega }}_0`$ be an isolated local strict minimum or maximum of $`\overline{\mathrm{\Gamma }}`$. There exists $`\overline{\epsilon }>0`$ such that if $`0<\epsilon <\overline{\epsilon }`$, then ($`\overline{๐ซ}_\epsilon `$) possesses a solution $`(u_\epsilon ,v_\epsilon )`$ such that $`u_\epsilon `$ concentrates in $`Q_\epsilon `$ with $`Q_\epsilon Q_0`$, as $`\epsilon 0`$, and $`v_\epsilon `$ concentrates in $`Q_\epsilon ^{}`$ with $`Q_\epsilon ^{}Q_0`$, as $`\epsilon 0`$.
###### Remark 1.6.
Let us observe that, if $`p=2`$ and $`N=3`$, then Theorem 1.1 is nothing else than a particular case of Theorem 1.5. Nevertheless, since problem ($`๐ซ_\epsilon `$) is more natural and more important by a physical point of view, we prefer to present Theorem 1.1 as our main result and to prove it directly, showing how, with slight modifications, the proof of Theorem 1.5 follows.
Theorem 1.1 will be proved as a particular case of a multiplicity result in Section 5 (see Theorem 5.1). The proof of the theorem relies on a finite dimensional reduction, precisely on the perturbation technique developed in . In Section 2 we give some preliminary lemmas and some estimates which will be useful in Section 3 and Section 4, where we perform the Liapunov-Schmidt reduction, making also the asymptotic expansion of the finite dimensional functional. Finally, in Section 5, we give also a short proof of Theorem 1.5.
Notation
* We denote $`\mathrm{\Omega }_0:=\mathrm{\Omega }(\mathrm{\Omega }\epsilon _0e_1)`$, where $`e_1=(1,0,0)`$ and $`\epsilon _0`$ is sufficiently small such that $`\mathrm{\Omega }_0\mathrm{}`$.
* If $`r>0`$ and $`x_0^3`$, $`B_r(x_0):=\{x^3:|xx_0|<r\}`$. We denote with $`B_r`$ the ball of radius $`r`$ centered in the origin.
* If $`u:^3`$ and $`P^3`$, we set $`u_P:=u(P)`$.
* If $`\epsilon >0`$, we set $`\mathrm{\Omega }_\epsilon :=\mathrm{\Omega }/\epsilon =\{x^3:\epsilon x\mathrm{\Omega }\}`$.
* We denote $`_\epsilon =H_0^1(\mathrm{\Omega }_\epsilon )\times H_0^1(\mathrm{\Omega }_\epsilon )`$.
* If there is no misunderstanding, we denote with $``$ and with $`()`$ respectively the norm and the scalar product both of $`H_0^1(\mathrm{\Omega }_\epsilon )`$ and of $`_\epsilon `$. While we denote with $`_^3`$ and with $`()_^3`$ respectively the norm and the scalar product of $`H^1(^3)`$.
* With $`C_i`$ and $`c_i`$, we denote generic positive constants, which may also vary from line to line.
## 2. Some preliminary
Performing the change of variable $`x\epsilon x`$, problem ($`๐ซ_\epsilon `$) becomes:
(2.1)
$$\{\begin{array}{cc}\mathrm{\Delta }u+J_1(\epsilon x)u=J_2(\epsilon x)u^3+\beta uv^2=0\hfill & \mathrm{in}\mathrm{\Omega }_\epsilon ,\hfill \\ \mathrm{\Delta }v+K_1(\epsilon x)v=K_2(\epsilon x)v^3+\beta u^2v=0\hfill & \mathrm{in}\mathrm{\Omega }_\epsilon ,\hfill \\ u,v>0\hfill & \mathrm{in}\mathrm{\Omega }_\epsilon ,\hfill \\ u=v=0\hfill & \mathrm{on}\mathrm{\Omega }_\epsilon ,\hfill \end{array}$$
where $`\mathrm{\Omega }_\epsilon =\epsilon ^1\mathrm{\Omega }`$. Of course if $`(u,v)`$ is a solution of (2.1), then $`(u(/\epsilon ),v(/\epsilon ))`$ is a solution of ($`๐ซ_\epsilon `$).
Solutions of (2.1) will be found in
$$_\epsilon =H_0^1(\mathrm{\Omega }_\epsilon )\times H_0^1(\mathrm{\Omega }_\epsilon ),$$
endowed with the following norm:
$$(u,v)__\epsilon ^2=u_{H_0^1(\mathrm{\Omega }_\epsilon )}^2+v_{H_0^1(\mathrm{\Omega }_\epsilon )}^2,\text{ for all }(u,v)_\epsilon .$$
If there is no misunderstanding, we denote with $``$ and with $`()`$ respectively the norm and the scalar product both of $`H^1(\mathrm{\Omega }_\epsilon )`$ and of $`_\epsilon `$.
Solutions of (2.1) are critical points of the functional $`f_\epsilon :_\epsilon `$, defined as
$`f_\epsilon (u,v)`$ $`={\displaystyle \frac{1}{2}}{\displaystyle _{\mathrm{\Omega }_\epsilon }}|u|^2+{\displaystyle \frac{1}{2}}{\displaystyle _{\mathrm{\Omega }_\epsilon }}J_1(\epsilon x)u^2{\displaystyle \frac{1}{4}}{\displaystyle _{\mathrm{\Omega }_\epsilon }}J_2(\epsilon x)u^4`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle _{\mathrm{\Omega }_\epsilon }}|v|^2+{\displaystyle \frac{1}{2}}{\displaystyle _{\mathrm{\Omega }_\epsilon }}K_1(\epsilon x)v^2{\displaystyle \frac{1}{4}}{\displaystyle _{\mathrm{\Omega }_\epsilon }}K_2(\epsilon x)v^4`$
$`{\displaystyle \frac{\beta }{2}}{\displaystyle _{\mathrm{\Omega }_\epsilon }}u^2v^2.`$
If we define $`f_\epsilon ^J:H_0^1(\mathrm{\Omega }_\epsilon )`$ and $`f_\epsilon ^K:H_0^1(\mathrm{\Omega }_\epsilon )`$ as
$`f_\epsilon ^J(u)=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _{\mathrm{\Omega }_\epsilon }}|u|^2+{\displaystyle \frac{1}{2}}{\displaystyle _{\mathrm{\Omega }_\epsilon }}J_1(\epsilon x)u^2{\displaystyle \frac{1}{4}}{\displaystyle _{\mathrm{\Omega }_\epsilon }}J_2(\epsilon x)u^4,`$
$`f_\epsilon ^K(v)=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _{\mathrm{\Omega }_\epsilon }}|v|^2+{\displaystyle \frac{1}{2}}{\displaystyle _{\mathrm{\Omega }_\epsilon }}K_1(\epsilon x)v^2{\displaystyle \frac{1}{4}}{\displaystyle _{\mathrm{\Omega }_\epsilon }}K_2(\epsilon x)v^4,`$
we have
$$f_\epsilon (u,v)=f_\epsilon ^J(u)+f_\epsilon ^K(v)\frac{\beta }{2}_{\mathrm{\Omega }_\epsilon }u^2v^2.$$
Furthermore, for any fixed $`Q\mathrm{\Omega }`$, we define the two functionals $`F_Q^J:H^1(^3)`$ and $`F_Q^K:H^1(^3)`$, as follows
$`F_Q^J(u)=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _^3}|u|^2+{\displaystyle \frac{1}{2}}{\displaystyle _^3}J_1(Q)u^2{\displaystyle \frac{1}{4}}{\displaystyle _^3}J_2(Q)u^4,`$
$`F_Q^K(v)=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _^3}|v|^2+{\displaystyle \frac{1}{2}}{\displaystyle _^3}K_1(Q)v^2{\displaystyle \frac{1}{4}}{\displaystyle _^3}K_2(Q)v^4.`$
The solutions of (2.1) will be found near $`(U^Q,V^Q)`$, properly truncated, where $`U^Q`$ is the unique solution of
(2.2)
$$\{\begin{array}{cc}\mathrm{\Delta }u+J_1(Q)u=J_2(Q)u^3\hfill & \text{in }^3,\hfill \\ u>0\hfill & \text{in }^3,\hfill \\ u(0)=\mathrm{max}_^3u,\hfill & \end{array}$$
and $`V_Q`$ is the unique solution of
(2.3)
$$\{\begin{array}{cc}\mathrm{\Delta }v+K_1(Q)v=K_2(Q)v^3\hfill & \text{in }^3,\hfill \\ v>0\hfill & \text{in }^3,\hfill \\ v(0)=\mathrm{max}_^3v,\hfill & \end{array}$$
for an appropriate choice of $`Q\mathrm{\Omega }_0`$. It is easy to see that
(2.4) $`U^Q(x)`$ $`=\sqrt{J_1(Q)/J_2(Q)}W\left(\sqrt{J_1(Q)}x\right),`$
(2.5) $`V^Q(x)`$ $`=\sqrt{K_1(Q)/K_2(Q)}W\left(\sqrt{K_1(Q)}x\right),`$
where $`W`$ is the unique solution of
(2.6)
$$\{\begin{array}{ccc}\mathrm{\Delta }z+z=z^3\hfill & \text{in }^3,\hfill & \\ z>0\hfill & \text{in }^3,\hfill & \\ z(0)=\mathrm{max}_^3z,\hfill & & \end{array}$$
which is radially symmetric and decays exponentially at infinity with its first derivatives (cf. ).
For all $`Q\mathrm{\Omega }_0`$, we set $`Q^{}=Q^{}(\epsilon ,Q)=Q+\sqrt{\epsilon }e_1\mathrm{\Omega }`$ and moreover we call $`P=P(\epsilon ,Q)=Q/\epsilon \mathrm{\Omega }_\epsilon `$ and $`P^{}=P^{}(\epsilon ,Q)=Q^{}/\epsilon \mathrm{\Omega }_\epsilon `$. Let us observe that
(2.7)
$$|PP^{}|=\frac{1}{\sqrt{\epsilon }}0,\mathrm{as}\epsilon 0.$$
Let $`\chi :^3`$ be a smooth function such that
(2.8)
$$\begin{array}{cccc}\hfill \chi (x)& =& 1,\hfill & \text{ for }|x|\epsilon ^{1/4};\hfill \\ \hfill \chi (x)& =& 0,\hfill & \text{ for }|x|2\epsilon ^{1/4};\hfill \\ \hfill 0\chi (x)& & 1,\hfill & \text{ for }\epsilon ^{1/4}|x|2\epsilon ^{1/4};\hfill \\ \hfill |\chi (x)|& & 2\epsilon ^{1/4},\hfill & \text{ for }\epsilon ^{1/4}|x|2\epsilon ^{1/4}.\hfill \end{array}$$
We denote
(2.9) $`U_P(x):=`$ $`\chi (xP)U^Q(xP),`$
(2.10) $`V_P^{}(x):=`$ $`\chi (xP^{})V^Q(xP^{}).`$
Let us observe that $`(U_P,V_P^{})_\epsilon `$. For $`Q`$ varying in $`\mathrm{\Omega }_0`$, $`(U_P,V_P^{})`$ describes a $`3`$-dimensional manifold, namely,
(2.11)
$$Z^\epsilon =\{(U_P,V_P^{}):Q\mathrm{\Omega }_0\}.$$
###### Remark 2.1.
Of course, if $`\mathrm{\Omega }=^3`$, then $`\mathrm{\Omega }_0=^3`$ and we do not need to truncate $`U^Q`$ and $`V^Q`$. In this case, we would have simply $`U_P=U^Q(P)`$ and $`V_P^{}=V^Q(P^{})`$.
First of all let us give the following estimate which will be very useful in the sequel.
###### Lemma 2.2.
For all $`Q\mathrm{\Omega }_0`$ and for all $`\epsilon `$ sufficiently small, if $`Q^{}=Q+\sqrt{\epsilon }e_1`$, $`P=Q/\epsilon \mathrm{\Omega }_\epsilon `$ and $`P^{}=Q^{}/\epsilon \mathrm{\Omega }_\epsilon `$, then
(2.12)
$$_{\mathrm{\Omega }_\epsilon }U_P^2V_P^{}^2=o(\epsilon ).$$
Proof Let us start observing that, since
$$|PP^{}|=\epsilon ^{1/2}>4\epsilon ^{1/4},$$
we infer that
$$B_{2\epsilon ^{1/4}}(P)B_{2\epsilon ^{1/4}}(P^{})=\mathrm{}.$$
Therefore, by the definitions of (2.9) and (2.10) and by the exponential decay of $`U_P`$ and $`V_P^{}`$, we get
$`{\displaystyle _{\mathrm{\Omega }_\epsilon }}U_P^2V_P^{}^2`$ $`{\displaystyle _{B_{2\epsilon ^{1/4}}(P)B_{2\epsilon ^{1/4}}(P^{})}}\left(U^Q\right)^2(xP)\left(V^Q\right)^2(xP^{})`$
$``$ $`c_1{\displaystyle _{^3B_{2\epsilon ^{1/4}}(P^{})}}\left(V^Q\right)^2(xP^{})`$
$`+c_2{\displaystyle _{^3B_{2\epsilon ^{1/4}}(P)}}\left(U^Q\right)^2(xP)=o(\epsilon ).`$
This concludes the proof. $`\mathrm{}`$
In the next lemma we show that the $`3`$-dimensional manifold $`Z_\epsilon `$, defined in (2.11), is actually a manifold of almost critical points of $`f_\epsilon `$.
###### Lemma 2.3.
For all $`Q\mathrm{\Omega }_0`$ and for all $`\epsilon `$ sufficiently small, if $`Q^{}=Q+\sqrt{\epsilon }e_1`$, $`P=Q/\epsilon \mathrm{\Omega }_\epsilon `$ and $`P^{}=Q^{}/\epsilon \mathrm{\Omega }_\epsilon `$, then
(2.13)
$$f_\epsilon (U_P,V_P^{})=O(\epsilon ^{1/2}).$$
Proof For all $`(u,v)_\epsilon `$, we have:
$`(f_\epsilon (U_P,V_P^{})(u,v))=`$ $`{\displaystyle _{\mathrm{\Omega }_\epsilon }}\left[U_Pu+J_1(\epsilon x)U_PuJ_2(\epsilon x)U_P^3u\right]`$
$`+{\displaystyle _{\mathrm{\Omega }_\epsilon }}\left[V_P^{}v+K_1(\epsilon x)V_P^{}vK_2(\epsilon x)V_P^{}^3v\right]`$
(2.14) $`\beta {\displaystyle _{\mathrm{\Omega }_\epsilon }}U_PV_P^{}^2u\beta {\displaystyle _{\mathrm{\Omega }_\epsilon }}U_P^2V_P^{}v.`$
Let us study the first integral of the right hand side of (2.14). By the exponential decay of $`U^Q`$ and recalling that $`U^Q`$ is solution of (2.2), we get
$`{\displaystyle _{\mathrm{\Omega }_\epsilon }}[U_P`$ $`u+J_1(\epsilon x)U_PuJ_2(\epsilon x)U_P^3u]`$
$`=`$ $`{\displaystyle _{(\mathrm{\Omega }Q)/\epsilon B_{\epsilon ^{1/4}}}}\left[U^Qu_P+J_1(\epsilon x+Q)U^Qu_P\right]`$
$`{\displaystyle _{(\mathrm{\Omega }Q)/\epsilon B_{\epsilon ^{1/4}}}}J_2(\epsilon x+Q)\left(U^Q\right)^3u_P+o(\epsilon )`$
$`=`$ $`{\displaystyle _^3}\left[U^Qu_P+J_1(\epsilon x+Q)U^Qu_PJ_2(\epsilon x+Q)\left(U^Q\right)^3u_P\right]+o(\epsilon )`$
$`=`$ $`{\displaystyle _^3}\left[U^Qu_P+J_1(Q)U^Qu_PJ_2(Q)\left(U^Q\right)^3u_P\right]`$
$`+{\displaystyle _^3}\left(J_1(\epsilon x+Q)J_1(Q)\right)U^Qu_P`$
$`{\displaystyle _^3}\left(J_2(\epsilon x+Q)J_2(Q)\right)\left(U^Q\right)^3u_P+o(\epsilon )`$
$`=`$ $`{\displaystyle _^3}\left(J_1(\epsilon x+Q)J_1(Q)\right)U^Qu_P`$
(2.15) $`{\displaystyle _^3}\left(J_2(\epsilon x+Q)J_2(Q)\right)\left(U^Q\right)^3u_P+o(\epsilon ).`$
Moreover, from the assumption $`DJ_i`$ bounded, we infer that
$$|J_i(\epsilon x+Q)J_i(Q)|c_1\epsilon |x|,$$
and so,
$`{\displaystyle _^3}\left(J_1(\epsilon x+Q)J_1(Q)\right)U^Qu_P`$ $`u\left({\displaystyle _^3}|J_1(\epsilon x+Q)J_1(Q)|^2|U^Q|^2\right)^{1/2}`$
(2.16) $``$ $`c_1u\left({\displaystyle _^3}\epsilon ^2|x|^2|U^Q|^2\right)^{1/2}=O(\epsilon )u.`$
Analogously,
(2.17)
$$_^3\left(J_2(\epsilon x+Q)J_2(Q)\right)\left(U^Q\right)^3u_P=O(\epsilon )u.$$
Therefore, by (2.15), (2.16) and (2.17), we infer
(2.18)
$$_{\mathrm{\Omega }_\epsilon }\left[U_Pu+J_1(\epsilon x)U_PuJ_2(\epsilon x)U_P^3u\right]=O(\epsilon )u.$$
Similarly, since $`V^Q`$ is solution of (2.3), we get
$`{\displaystyle _{\mathrm{\Omega }_\epsilon }}[V_P^{}v`$ $`+K_1(\epsilon x)V_P^{}vK_2(\epsilon x)V_P^{}^3v]`$
$`=`$ $`{\displaystyle _^3}\left(K_1(\epsilon x+Q+\sqrt{\epsilon }e_1)K_1(Q)\right)V^Qv_P^{}`$
(2.19) $`{\displaystyle _^3}\left(K_2(\epsilon x+Q+\sqrt{\epsilon }e_1)K_2(Q)\right)\left(V^Q\right)^3v_P^{}+o(\epsilon ).`$
Therefore, from the assumption $`DK_i`$ bounded, we infer that
$$|K_i(\epsilon x+Q+\sqrt{\epsilon }e_1)K_i(Q)|c_2\sqrt{\epsilon }|\sqrt{\epsilon }x+e_1|,$$
and so,
$`{\displaystyle _^3}(K_1(\epsilon x+Q+`$ $`\sqrt{\epsilon }e_1)K_1(Q))V^Qv_P^{}`$
$`v\left({\displaystyle _^3}|K_1(\epsilon x+Q+\sqrt{\epsilon }e_1)K_1(Q)|^2|V^Q|^2\right)^{1/2}`$
(2.20) $`c_2v\left({\displaystyle _^3}\epsilon |\sqrt{\epsilon }x+e_1|^2|V^Q|^2\right)^{1/2}=O(\epsilon ^{1/2})v.`$
Analogously,
(2.21)
$$_^3\left(K_2(\epsilon x+Q+\sqrt{\epsilon }e_1)K_2(Q)\right)\left(V^Q\right)^3v_P^{}=O(\epsilon ^{1/2})v.$$
Therefore, by (2.19), (2.20) and (2.21), we infer
(2.22)
$$_{\mathrm{\Omega }_\epsilon }\left[V_P^{}v+K_1(\epsilon x)V_P^{}vK_2(\epsilon x)V_P^{}^3v\right]=O(\epsilon ^{1/2})v.$$
Let us study the last two terms of (2.14). Arguing as in Lemma 2.2, we get
(2.23)
$$\left|_{\mathrm{\Omega }_\epsilon }U_PV_P^{}^2u\right|c_3\left(_{\mathrm{\Omega }_\epsilon }U_P^{4/3}V_P^{}^{8/3}\right)^{3/4}u=o(\epsilon )u,$$
and
(2.24)
$$\left|_{\mathrm{\Omega }_\epsilon }U_P^2V_P^{}v\right|=o(\epsilon )v.$$
Now the conclusion of the proof easily follows by (2.14), (2.18), (2.22), (2.23) and (2.24). $`\mathrm{}`$
## 3. Invertibility of $`D^2f_\epsilon `$ on $`\left(T_{(U_P,V_P^{})}Z^\epsilon \right)^{}`$
In this section we will show that $`D^2f_\epsilon `$ is invertible on $`\left(T_{(U_P,V_P^{})}Z^\epsilon \right)^{}`$, where $`T_{(U_P,V_P^{})}Z^\epsilon `$ denotes the tangent space to $`Z^\epsilon `$ at the point $`(U_P,V_P^{})`$.
Let $`L_{\epsilon ,Q}:(T_{(U_P,V_P^{})}Z^\epsilon )^{}(T_{(U_P,V_P^{})}Z^\epsilon )^{}`$ denote the operator defined by setting $`(L_{\epsilon ,Q}(h,h^{})(k,k^{}))=D^2f_\epsilon (U_P,V_P^{})[(h,h^{}),(k,k^{})]`$.
###### Lemma 3.1.
Given $`\mu >0`$, there exists $`C>0`$ such that, for $`\epsilon `$ small enough and for all $`Q\mathrm{\Omega }_0`$ with $`|Q|\mu `$, one has that
(3.1)
$$L_{\epsilon ,Q}(h,h^{})C(h,h^{}),(h,h^{})(T_{(U_P,V_P^{})}Z^\epsilon )^{}.$$
Proof First of all, let us observe that, for all $`(h,h^{}),(k,k^{})_\epsilon `$, we have
(3.2)
$$\begin{array}{c}D^2f_\epsilon (u,v)[(h,h^{}),(k,k^{})]=D^2f_\epsilon ^J(u)[h,k]+D^2f_\epsilon ^K(v)[h^{},k^{}]\hfill \\ \hfill \beta _{\mathrm{\Omega }_\epsilon }v^2hk2\beta _{\mathrm{\Omega }_\epsilon }uvhk^{}2\beta _{\mathrm{\Omega }_\epsilon }uvh^{}k\beta _{\mathrm{\Omega }_\epsilon }u^2h^{}k^{}.\end{array}$$
By (2.4), if we set $`a(Q)=\sqrt{J_1(Q)/J_2(Q)}`$ and $`b(Q)=\sqrt{J_1(Q)}`$, we have that $`U^Q(x)=a(Q)W(b(Q)x)`$ and so $`U_P(x)=\chi (xP)a(\epsilon P)W(b(\epsilon P)(xP))`$. Therefore, we have:
$`_{P_i}U_P(x)=`$ $`_{P_i}\left(\chi (xP)U^Q(xP)\right)`$
$`=`$ $`U^Q(xP)_{x_i}\chi (xP)+\chi (xP)_{P_i}U^Q(xP)`$
$`=`$ $`U^Q(xP)_{x_i}\chi (xP)+\epsilon \chi (xP)_{P_i}a(\epsilon P)W(b(\epsilon P)(xP))`$
$`+\epsilon \chi (xP)a(\epsilon P)_{P_i}a(\epsilon P)W(b(\epsilon P)(xP))(xP)`$
$`\chi (xP)a(\epsilon P)b(\epsilon P)(_{x_i}W)(b(\epsilon P)(xP)).`$
Hence
(3.3)
$$_{P_i}U_P(x)=_{x_i}U_P(x)+O(\epsilon ).$$
Analogously, we can prove that
(3.4)
$$_{P_i}V_P^{}(x)=_{P_i^{}}V_P^{}(x)=_{x_i}V_P^{}(x)+O(\epsilon ).$$
We recall that
$$T_{(U_P,V_P^{})}Z^\epsilon =\mathrm{span}__\epsilon \{(_{P_1}U_P,_{P_1}V_P^{}),(_{P_2}U_P,_{P_2}V_P^{}),(_{P_3}U_P,_{P_3}V_P^{})\}.$$
We set
$$๐ฑ_\epsilon =\mathrm{span}__\epsilon \{(U_P,V_P^{}),(_{x_1}U_P,_{x_1}V_P^{}),(_{x_2}U_P,_{x_2}V_P^{}),(_{x_3}U_P,_{x_3}V_P^{})\}.$$
By (3.3) and (3.4), therefore it suffices to prove equation (3.1) for all $`(h,h^{})`$ $`\mathrm{span}__\epsilon \{(U_P,V_P^{}),(\varphi ,\varphi ^{})\}`$, where $`(\varphi ,\varphi ^{})`$ is orthogonal to $`๐ฑ_\epsilon `$. Precisely we shall prove that there exist $`C_1,C_2>0`$ such that, for all $`\epsilon >0`$ small enough, one has:
(3.5) $`(L_{\epsilon ,Q}(U_P,V_P^{})(U_P,V_P^{}))`$ $`C_1<0,`$
(3.6) $`(L_{\epsilon ,Q}(\varphi ,\varphi ^{})(\varphi ,\varphi ^{}))`$ $`C_2(\varphi ,\varphi ^{})^2,\text{for all }(\varphi ,\varphi ^{})๐ฑ_\epsilon .`$
Proof of (3.5). By (3.2), we get:
(3.7)
$$\begin{array}{c}D^2f_\epsilon (U_P,V_P^{})[(U_P,V_P^{}),(U_P,V_P^{})]\hfill \\ \hfill =D^2f_\epsilon ^J(U_P)[U_P,U_P]+D^2f_\epsilon ^K(V_P^{})[V_P^{},V_P^{}]6\beta _{\mathrm{\Omega }_\epsilon }U_P^2V_P^{}^2.\end{array}$$
Let us study the first term of the right hand side of (3.7).
$`D^2f_\epsilon ^J(U_P)[U_P,U_P]=`$ $`{\displaystyle _{\mathrm{\Omega }_\epsilon }}|U_P|^2+{\displaystyle _{\mathrm{\Omega }_\epsilon }}J_1(\epsilon x)U_P^23{\displaystyle _{\mathrm{\Omega }_\epsilon }}J_2(\epsilon x)U_P^4`$
$`=`$ $`{\displaystyle \underset{(\mathrm{\Omega }Q)/\epsilon B_{\epsilon ^{1/4}}}{}}\left[|U^Q|^2+J_1(\epsilon x+Q)\left(U^Q\right)^23J_2(\epsilon x+Q)\left(U^Q\right)^4\right]`$
$`+o(\epsilon )`$
$`=`$ $`{\displaystyle _^3}\left[|U^Q|^2+J_1(Q)\left(U^Q\right)^23J_2(Q)\left(U^Q\right)^4\right]`$
$`+{\displaystyle _^3}\left(J_1(\epsilon x+Q)J_1(Q)\right)\left(U^Q\right)^2`$
$`3{\displaystyle _^3}\left(J_2(\epsilon x+Q)J_2(Q)\right)\left(U^Q\right)^4+o(\epsilon )`$
$`=`$ $`2{\displaystyle _^3}J_2(Q)\left(U^Q\right)^4+O(\epsilon )`$
$`=`$ $`2J_1(Q)^{\frac{1}{2}}J_2(Q)^1{\displaystyle _^3}W^4+O(\epsilon )c_1.`$
In a similar way it is possible to prove that
$$D^2f_\epsilon ^K(V_P^{})[V_P^{},V_P^{}]c_2.$$
Finally, by Lemma 2.2, we know that
$$_{\mathrm{\Omega }_\epsilon }U_P^2V_P^{}^2=o(\epsilon ),$$
and so equation (3.5) is proved.
Proof of (3.6). Recalling the definition of $`\chi `$, (see (2.8)), we set $`\chi _1:=\chi `$ and $`\chi _2:=1\chi _1`$. Given $`(\varphi ,\varphi ^{})๐ฑ_\epsilon `$, let us consider the functions
(3.8) $`\varphi _i(x)=\chi _i(xP)\varphi (x),i=1,2;`$
(3.9) $`\varphi _i^{}(x)=\chi _i(xP^{})\varphi ^{}(x),i=1,2.`$
With calculations similar to those of , we have
(3.10) $`\varphi ^2=`$ $`\varphi _1^2+\varphi _2^2+\underset{I_\varphi }{\underset{}{2{\displaystyle _{\mathrm{\Omega }_\epsilon }}\chi _1\chi _2(\varphi ^2+|\varphi |^2)}}+O(\epsilon ^{1/4})\varphi ^2,`$
(3.11) $`\varphi ^{}^2=`$ $`\varphi _1^{}^2+\varphi _2^{}^2+\underset{I_\varphi ^{}}{\underset{}{2{\displaystyle _{\mathrm{\Omega }_\epsilon }}\chi _1\chi _2((\varphi ^{})^2+|\varphi ^{}|^2)}}+O(\epsilon ^{1/4})\varphi ^{}^2.`$
We need to evaluate the three terms in the equation below:
$`(L_{\epsilon ,Q}(\varphi ,\varphi ^{})(\varphi ,\varphi ^{}))=`$ $`(L_{\epsilon ,Q}(\varphi _1,\varphi _1^{})(\varphi _1,\varphi _1^{}))+(L_{\epsilon ,Q}(\varphi _2,\varphi _2^{})(\varphi _2,\varphi _2^{}))`$
(3.12) $`+2(L_{\epsilon ,Q}(\varphi _1,\varphi _1^{})(\varphi _2,\varphi _2^{})).`$
Let us start with $`(L_{\epsilon ,Q}(\varphi _1,\varphi _1^{})(\varphi _1,\varphi _1^{}))`$. Since $`\beta <0`$, we get
$`(L_{\epsilon ,Q}(\varphi _1,\varphi _1^{})(\varphi _1,\varphi _1^{}))=`$ $`D^2f_\epsilon ^J(U_P)[\varphi _1,\varphi _1]+D^2f_\epsilon ^K(V_P^{})[\varphi _1^{},\varphi _1^{}]`$
$`4\beta {\displaystyle _{\mathrm{\Omega }_\epsilon }}U_PV_P^{}\varphi _1\varphi _1^{}\beta {\displaystyle _{\mathrm{\Omega }_\epsilon }}U_P^2\varphi _{1}^{}{}_{}{}^{2}\beta {\displaystyle _{\mathrm{\Omega }_\epsilon }}V_P^{}^2\varphi _1^2`$
$`>`$ $`D^2f_\epsilon ^J(U_P)[\varphi _1,\varphi _1]+D^2f_\epsilon ^K(V_P^{})[\varphi _1^{},\varphi _1^{}]`$
(3.13) $`4\beta {\displaystyle _{\mathrm{\Omega }_\epsilon }}U_PV_P^{}\varphi _1\varphi _1^{}.`$
Arguing as in Lemma 2.2, we know that
(3.14)
$$_{\mathrm{\Omega }_\epsilon }U_PV_P^{}\varphi _1\varphi _1^{}=o(\epsilon ).$$
Therefore we need only to study the first two terms of the right hand side of (3.13). For simplicity, we can assume that $`Q=\epsilon P`$ is the origin $`๐ช`$. In this case, we recall that we denote with $`U^๐ช`$ the unique solution of (2.2) whenever $`Q=๐ช`$, while we denote with $`U_๐ช`$ the truncation of $`U^๐ช`$, namely $`U_๐ช=\chi U^๐ช`$, where $`\chi `$ is defined in (2.8). We have
$`D^2f_\epsilon ^J(U_๐ช)[\varphi _1,\varphi _1]=`$ $`{\displaystyle _{\mathrm{\Omega }_\epsilon }}\left[|\varphi _1|^2+J_1(\epsilon x)\varphi _1^23J_2(\epsilon x)U_๐ช^2\varphi _1^2\right]`$
$`=`$ $`{\displaystyle _^3}\left[|\varphi _1|^2+J_1(\epsilon x)\varphi _1^23J_2(\epsilon x)\left(U^๐ช\right)^2\varphi _1^2\right]+o(\epsilon )\varphi ^2`$
$`=`$ $`D^2F^{J(๐ช)}(U^๐ช)[\varphi _1,\varphi _1]`$
$`+{\displaystyle _^3}\left(J_1(\epsilon x)J_1(๐ช)\right)\varphi _1^2`$
$`3{\displaystyle _^3}\left(J_2(\epsilon x)J_2(๐ช)\right)\left(U^๐ช\right)^2\varphi _1^2+o(\epsilon )\varphi ^2`$
$``$ $`D^2F^{J(๐ช)}(U^๐ช)[\varphi _1,\varphi _1]c_3\epsilon {\displaystyle _^3}|x|\varphi _1^2+O(\epsilon )\varphi ^2`$
$`=`$ $`D^2F^{J(๐ช)}(U^๐ช)[\varphi _1,\varphi _1]+O(\epsilon ^{3/4})\varphi ^2,`$
therefore
(3.15)
$$D^2f_\epsilon ^J(U_๐ช)[\varphi _1,\varphi _1]D^2F^{J(๐ช)}(U^๐ช)[\varphi _1,\varphi _1]+O(\epsilon ^{3/4})\varphi ^2.$$
We recall that $`\varphi `$ is orthogonal to
$$๐ฑ_\epsilon ^U=\mathrm{span}_{H_0^1(\mathrm{\Omega }_\epsilon )}\{U_๐ช,_{x_1}U_๐ช,_{x_2}U_๐ช,_{x_3}U_๐ช\}.$$
Moreover by , we know that if $`\stackrel{~}{\varphi }`$ is orthogonal to $`๐ฑ`$ with
$$๐ฑ^U=\mathrm{span}_{H^1(^3)}\{U^๐ช,_{x_1}U^๐ช,_{x_2}U^๐ช,_{x_3}U^๐ช\},$$
then the fact that $`U^๐ช`$ is a Mountain Pass critical point of $`F^{J(๐ช)}`$ implies that
(3.16)
$$D^2F^{J(๐ช)}(U^๐ช)[\stackrel{~}{\varphi },\stackrel{~}{\varphi }]>c_4\stackrel{~}{\varphi }_^3^2\text{for all }\stackrel{~}{\varphi }๐ฑ^U.$$
We can write $`\varphi _1=\xi +\zeta `$, where $`\xi ๐ฑ^U`$ and $`\zeta ๐ฑ^U`$. More precisely
$$\xi =(\varphi _1U^๐ช)_^3U^๐ชU^๐ช_^3^2+\underset{i=1}{\overset{3}{}}(\varphi _1_{x_i}U^๐ช)_^3_{x_i}U^๐ช_{x_i}U^๐ช_^3^2.$$
Let us calculate $`(\varphi _1U^๐ช)_^3`$. By the exponential decay of $`U^๐ช`$ and since $`\varphi ๐ฑ_\epsilon ^U`$, we have
$`(\varphi _1U^๐ช)_^3=`$ $`{\displaystyle _^3}\varphi _1U^๐ช+{\displaystyle _^3}\varphi _1U^๐ช`$
$`=`$ $`{\displaystyle _{\mathrm{\Omega }_\epsilon }}\varphi _1U_๐ช+{\displaystyle _{\mathrm{\Omega }_\epsilon }}\varphi _1U_๐ช+o(\epsilon )\varphi `$
$`=`$ $`{\displaystyle _{\mathrm{\Omega }_\epsilon }}\varphi U_๐ช+{\displaystyle _{\mathrm{\Omega }_\epsilon }}\varphi U_๐ช+o(\epsilon )\varphi =o(\epsilon )\varphi .`$
In a similar way, we can prove also that $`(\varphi _1_{x_i}U^๐ช)_^3=o(\epsilon )\varphi `$, and so
(3.17) $`\xi _^3=o(\epsilon )\varphi ,`$
(3.18) $`\zeta _^3=\varphi _1+o(\epsilon )\varphi .`$
Let us estimate $`D^2F^{J(๐ช)}(U^๐ช)[\varphi _1,\varphi _1]`$. We get:
$`D^2F^{J(๐ช)}(U^๐ช)[\varphi _1,\varphi _1]=`$ $`D^2F^{J(๐ช)}(U^๐ช)[\zeta ,\zeta ]+2D^2F^{J(๐ช)}(U^๐ช)[\zeta ,\xi ]`$
(3.19) $`+D^2F^{J(๐ช)}(U^๐ช)[\xi ,\xi ].`$
By (3.16) and (3.18), since $`\zeta ๐ฑ^U`$, we know that
$$D^2F^{J(๐ช)}(U^๐ช)[\zeta ,\zeta ]>c_3\zeta _^3^2=c_3\varphi _1^2+o(\epsilon )\varphi ^2,$$
while, by (3.17) and straightforward calculations, we have
$`D^2F^{J(๐ช)}(U^๐ช)[\zeta ,\xi ]=o(\epsilon )\varphi ^2,`$
$`D^2F^{J(๐ช)}(U^๐ช)[\xi ,\xi ]=o(\epsilon )\varphi ^2.`$
By these last two estimates, (3.19) and (3.15), we can say that
$$D^2f_\epsilon ^J(U_๐ช)[\varphi _1,\varphi _1]>c_4\varphi _1^2+O(\epsilon ^{3/4})\varphi ^2.$$
Hence, in the general case, we infer that, for all $`Q\mathrm{\Omega }_0`$ with $`|Q|\mu `$,
(3.20)
$$D^2f_\epsilon ^J(U_P)[\varphi _1,\varphi _1]>c_4\varphi _1^2+O(\epsilon ^{3/4})\varphi ^2,$$
and, analogously,
(3.21)
$$D^2f_\epsilon ^K(V_P^{})[\varphi _1^{},\varphi _1^{}]>c_5\varphi _1^{}^2+O(\epsilon ^{1/2})\varphi ^{}^2.$$
By (3.13), (3.14), (3.20) and (3.21), we can say that
(3.22)
$$(L_{\epsilon ,Q}(\varphi _1,\varphi _1^{})(\varphi _1,\varphi _1^{}))>c_6(\varphi _1,\varphi _1^{})^2+O(\epsilon ^{1/2})(\varphi ,\varphi _1^{})^2.$$
Let us now evaluate $`(L_{\epsilon ,Q}(\varphi _2,\varphi _2^{})(\varphi _2,\varphi _2^{}))`$. Arguing as in Lemma 2.2, since $`\beta <0`$ and using the definition of $`\chi _i`$ and the exponential decay of $`U_P`$ and of $`V_P^{}`$, we easily get:
$`(L_{\epsilon ,Q}(\varphi _2,\varphi _2^{})(\varphi _2,\varphi _2^{}))=`$ $`D^2f_\epsilon ^J(U_P)[\varphi _2,\varphi _2]+D^2f_\epsilon ^K(V_P^{})[\varphi _2^{},\varphi _2^{}]`$
$`4\beta {\displaystyle _{\mathrm{\Omega }_\epsilon }}U_PV_P^{}\varphi _2\varphi _2^{}\beta {\displaystyle _{\mathrm{\Omega }_\epsilon }}U_P^2\varphi _{2}^{}{}_{}{}^{2}\beta {\displaystyle _{\mathrm{\Omega }_\epsilon }}V_P^{}^2\varphi _2^2`$
$``$ $`D^2f_\epsilon ^J(U_P)[\varphi _2,\varphi _2]+D^2f_\epsilon ^K(V_P^{})[\varphi _2^{},\varphi _2^{}]+o(\epsilon )(\varphi ,\varphi ^{})^2`$
(3.23) $``$ $`c_7(\varphi _2,\varphi _2^{})^2+o(\epsilon )(\varphi ,\varphi ^{})^2.`$
Let us now study $`(L_{\epsilon ,Q}(\varphi _1,\varphi _1^{})(\varphi _2,\varphi _2^{}))`$. Arguing as in Lemma 2.2, we get
$`(L_{\epsilon ,Q}(\varphi _1,\varphi _1^{})(\varphi _2,\varphi _2^{}))=`$ $`D^2f_\epsilon ^J(U_P)[\varphi _1,\varphi _2]+D^2f_\epsilon ^K(V_P^{})[\varphi _1^{},\varphi _2^{}]`$
$`2\beta {\displaystyle _{\mathrm{\Omega }_\epsilon }}U_PV_P^{}\varphi _1\varphi _2^{}2\beta {\displaystyle _{\mathrm{\Omega }_\epsilon }}U_PV_P^{}\varphi _2\varphi _1^{}`$
$`\beta {\displaystyle _{\mathrm{\Omega }_\epsilon }}U_P^2\varphi _1^{}\varphi _2^{}\beta {\displaystyle _{\mathrm{\Omega }_\epsilon }}V_P^{}^2\varphi _1\varphi _2`$
$`=`$ $`D^2f_\epsilon ^J(U_P)[\varphi _1,\varphi _2]+D^2f_\epsilon ^K(V_P^{})[\varphi _1^{},\varphi _2^{}]`$
(3.24) $`\beta {\displaystyle _{\mathrm{\Omega }_\epsilon }}U_P^2\varphi _1^{}\varphi _2^{}\beta {\displaystyle _{\mathrm{\Omega }_\epsilon }}V_P^{}^2\varphi _1\varphi _2+o(\epsilon )(\varphi ,\varphi ^{})^2.`$
Using the definition of $`\chi _i`$ and the exponential decay of $`U_P`$ and of $`V_P^{}`$, we easily get:
(3.25) $`D^2f_\epsilon ^J(U_P)[\varphi _1,\varphi _2]c_8I_\varphi +O(\epsilon ^{1/4})\varphi ^2,`$
(3.26) $`D^2f_\epsilon ^K(V_P^{})[\varphi _1^{},\varphi _2^{}]c_9I_\varphi ^{}+O(\epsilon ^{1/4})\varphi ^{}^2,`$
where $`I_\varphi `$ and $`I_\varphi ^{}`$ are defined, respectively in (3.10) and (3.11). Moreover, by the definition of $`\chi `$, (see (2.8)), and by the definitions of $`\varphi _i`$ and $`\varphi _i^{}`$, (see (3.8) and (3.9)),
$$\varphi _1(x)\varphi _2(x)=\chi (xP)(1\chi (xP))\varphi ^2(x)0,\text{for all }x^3,$$
and so, also
$$\varphi _1^{}(x)\varphi _2^{}(x)0,\text{for all }x^3.$$
Therefore
(3.27)
$$\beta _{\mathrm{\Omega }_\epsilon }U_P^2\varphi _1^{}\varphi _2^{}\beta _{\mathrm{\Omega }_\epsilon }V_P^{}^2\varphi _1\varphi _20.$$
By (3.24), (3.25), (3.26) and (3.27), we infer
(3.28)
$$(L_{\epsilon ,Q}(\varphi _1,\varphi _1^{})(\varphi _2,\varphi _2^{}))c_{10}\left(I_\varphi +I_\varphi ^{}\right)+O(\epsilon ^{1/4})(\varphi ,\varphi ^{})^2.$$
Hence, by (3.12), (3.22), (3.23), (3.28) and recalling (3.10) and (3.11), we get
$$(L_{\epsilon ,Q}(\varphi ,\varphi ^{})(\varphi ,\varphi ^{}))c_{11}(\varphi ,\varphi ^{})^2+O(\epsilon ^{1/4})(\varphi ,\varphi ^{})^2.$$
This completes the proof of the lemma. $`\mathrm{}`$
## 4. The finite dimensional reduction
By means of the Liapunov-Schmidt reduction, the existence of critical points of $`f_\epsilon `$ can be reduced to the search of critical points of an auxiliary finite-dimensional functional.
###### Lemma 4.1.
Fix $`\mu >0`$. For $`\epsilon >0`$ small enough and for all $`Q\mathrm{\Omega }_0`$ with $`|Q|\mu `$, there exists a unique $`(w,w^{})=(w(\epsilon ,Q),w^{}(\epsilon ,Q))_\epsilon `$ of class $`C^1`$ such that:
1. $`(w(\epsilon ,Q),w^{}(\epsilon ,Q))(T_{(U_P,V_P^{})}Z^\epsilon )^{}`$;
2. $`f_\epsilon (U_P+w,V_P^{}+w^{})T_{(U_P,V_P^{})}Z^\epsilon `$.
Moreover, the functional $`๐_\epsilon :\mathrm{\Omega }_0`$, defined as:
$$๐_\epsilon (Q):=f_\epsilon (U_{Q/\epsilon }+w(\epsilon ,Q),V_{(Q+e_1\sqrt{\epsilon })/\epsilon }+w^{}(\epsilon ,Q))$$
is of class $`C^1`$ and satisfies:
$$๐_\epsilon (Q_0)=0f_\epsilon (U_{Q_0/\epsilon }+w(\epsilon ,Q_0),V_{(Q_0+e_1\sqrt{\epsilon })/\epsilon }+w^{}(\epsilon ,Q_0))=0.$$
Proof Let $`๐ซ=๐ซ_{\epsilon ,Q}`$ denote the projection onto $`(T_{(U_P,V_P^{})}Z^\epsilon )^{}`$. We want to find a solution $`(w,w^{})(T_{(U_P,V_P^{})}Z^\epsilon )^{}`$ of the equation
$$Pf_\epsilon (U_P+w,V_P^{}+w^{})=0.$$
One has that
$$f_\epsilon (U_P+w,V_P^{}+w^{})=f_\epsilon (U_P,V_P^{})+D^2f_\epsilon (U_P,V_P^{})[w,w^{}]+R(U_P,V_P^{},w,w^{})$$
with $`R(U_P,V_P^{},w,w^{})=o((w,w^{}))`$, uniformly with respect to $`(U_P,V_P^{})`$. Therefore, our equation is:
(4.1)
$$L_{\epsilon ,Q}(w,w^{})+๐ซf_\epsilon (U_P,V_P^{})+๐ซR(U_P,V_P^{},w,w^{})=0.$$
According to Lemma 3.1, this is equivalent to
$$(w,w^{})=N_{\epsilon ,Q}(w,w^{}),$$
where
$$N_{\epsilon ,Q}(w,w^{})=\left(L_{\epsilon ,Q}\right)^1\left(๐ซf_\epsilon (U_P,V_P^{})+๐ซR(U_P,V_P^{},w,w^{})\right).$$
By (2.13) it follows that
(4.2)
$$N_{\epsilon ,Q}(w,w^{})=O(\epsilon ^{1/2})+o((w,w^{})).$$
Therefore it is easy to check that $`N_{\epsilon ,Q}`$ is a contraction on some ball in $`(T_{(U_P,V_P^{})}Z^\epsilon )^{}`$ provided that $`\epsilon >0`$ is small enough. Then there exists a unique $`(w,w^{})`$ such that $`(w,w^{})=N_{\epsilon ,Q}(w,w^{})`$. Let us point out that we cannot use the Implicit Function Theorem to find $`(w(\epsilon ,Q),w^{}(\epsilon ,Q))`$, because the map $`(\epsilon ,u,v)๐ซf_\epsilon (u,v)`$ fails to be $`C^2`$. However, fixed $`\epsilon >0`$ small, we can apply the Implicit Function Theorem to the map $`(Q,w,w^{})๐ซf_\epsilon (U_P+w,V_P^{}+w^{})`$. Then, in particular, the function $`(w(\epsilon ,Q),w^{}(\epsilon ,Q))`$ turns out to be of class $`C^1`$ with respect to $`Q`$. Finally, it is a standard argument, see , to check that the critical points of $`๐_\epsilon (Q)=f_\epsilon (U_P+w,V_P^{}+w^{})`$ give rise to critical points of $`f_\epsilon `$. $`\mathrm{}`$
###### Remark 4.2.
From (4.2) it immediately follows that:
(4.3)
$$(w,w^{})=O(\epsilon ^{1/2}).$$
Let us now make the asymptotic expansion of the finite dimensional functional.
###### Theorem 4.3.
Fix $`\mu >0`$ and let $`Q\mathrm{\Omega }_0`$ with $`|Q|\mu `$, $`Q^{}=Q+\sqrt{\epsilon }e_1`$, $`P=Q/\epsilon \mathrm{\Omega }_\epsilon `$ and $`P^{}=Q^{}/\epsilon \mathrm{\Omega }_\epsilon `$. Suppose (J) and (K). Then, for $`\epsilon `$ sufficiently small, we get:
(4.4)
$$๐_\epsilon (Q)=f_\epsilon (U_P+w(\epsilon ,Q),V_P^{}+w^{}(\epsilon ,Q))=c_0\mathrm{\Gamma }(Q)+o(\epsilon ^{1/4}),$$
where $`\mathrm{\Gamma }:\mathrm{\Omega }_0`$ is defined in (1.4), namely
$$\mathrm{\Gamma }(Q)=J_1(Q)^{\frac{1}{2}}J_2(Q)^1+K_1(Q)^{\frac{1}{2}}K_2(Q)^1,$$
and
(4.5)
$$c_0:=\frac{1}{2}_^3W^4$$
with $`W`$ the unique solution of (2.6).
Proof We have:
$`๐_\epsilon (Q)=`$ $`f_\epsilon (U_P+w(\epsilon ,Q),V_P^{}+w^{}(\epsilon ,Q))`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _{\mathrm{\Omega }_\epsilon }}|(U_P+w)|^2+{\displaystyle \frac{1}{2}}{\displaystyle _{\mathrm{\Omega }_\epsilon }}J_1(\epsilon x)(U_P+w)^2{\displaystyle \frac{1}{4}}{\displaystyle _{\mathrm{\Omega }_\epsilon }}J_2(\epsilon x)(U_P+w)^4`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle _{\mathrm{\Omega }_\epsilon }}|(V_P^{}+w^{})|^2+{\displaystyle \frac{1}{2}}{\displaystyle _{\mathrm{\Omega }_\epsilon }}K_1(\epsilon x)(V_P^{}+w^{})^2{\displaystyle \frac{1}{4}}{\displaystyle _{\mathrm{\Omega }_\epsilon }}K_2(\epsilon x)(V_P^{}+w^{})^4`$
$`{\displaystyle \frac{\beta }{2}}{\displaystyle _{\mathrm{\Omega }_\epsilon }}(U_P+w)^2(V_P^{}+w^{})^2.`$
Therefore, by (4.3) and Lemma 2.2,
$`๐_\epsilon (Q)=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _{\mathrm{\Omega }_\epsilon }}|U_P|^2+{\displaystyle \frac{1}{2}}{\displaystyle _{\mathrm{\Omega }_\epsilon }}J_1(\epsilon x)U_P^2{\displaystyle \frac{1}{4}}{\displaystyle _{\mathrm{\Omega }_\epsilon }}J_2(\epsilon x)U_P^4`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle _{\mathrm{\Omega }_\epsilon }}|V_P^{}|^2+{\displaystyle \frac{1}{2}}{\displaystyle _{\mathrm{\Omega }_\epsilon }}K_1(\epsilon x)V_P^{}^2{\displaystyle \frac{1}{4}}{\displaystyle _{\mathrm{\Omega }_\epsilon }}K_2(\epsilon x)V_P^{}^4+O(\epsilon ^{1/2})`$
(4.6) $`=`$ $`f_\epsilon ^J(U_P)+f_\epsilon ^K(V_P^{})+O(\epsilon ^{1/2}).`$
Let us study the first term of the right hand side of (4.6).
$`f_\epsilon ^J(U_P)=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _{\mathrm{\Omega }_\epsilon }}|U_P|^2+{\displaystyle \frac{1}{2}}{\displaystyle _{\mathrm{\Omega }_\epsilon }}J_1(\epsilon x)U_P^2{\displaystyle \frac{1}{4}}{\displaystyle _{\mathrm{\Omega }_\epsilon }}J_2(\epsilon x)U_P^4`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{(\mathrm{\Omega }Q)/\epsilon B_{\epsilon ^{1/4}}}{}}\left[|U^Q|^2+J_1(\epsilon x+Q)\left(U^Q\right)^2{\displaystyle \frac{1}{2}}J_2(\epsilon x+Q)\left(U^Q\right)^4\right]+o(\epsilon )`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _^3}|U^Q|^2+{\displaystyle \frac{1}{2}}{\displaystyle _^3}J_1(Q)\left(U^Q\right)^2{\displaystyle \frac{1}{4}}{\displaystyle _^3}J_2(Q)\left(U^Q\right)^4`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle _^3}\left(J_1(\epsilon x+Q)J_1(Q)\right)\left(U^Q\right)^2`$
$`{\displaystyle \frac{1}{4}}{\displaystyle _^3}\left(J_2(\epsilon x+Q)J_2(Q)\right)\left(U^Q\right)^4+o(\epsilon )`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _^3}J_2(Q)\left(U^Q\right)^4+o(\epsilon ^{1/4})`$
$`=`$ $`{\displaystyle \frac{1}{2}}J_1(Q)^{\frac{1}{2}}J_2(Q)^1{\displaystyle _^3}W^4+o(\epsilon ^{1/4}).`$
Hence
(4.7)
$$f_\epsilon ^J(U_P)=\frac{1}{2}J_1(Q)^{\frac{1}{2}}J_2(Q)^1_^3W^4+o(\epsilon ^{1/4}).$$
Analogously,
$`f_\epsilon ^K(V_P^{})=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _{\mathrm{\Omega }_\epsilon }}|V_P^{}|^2+{\displaystyle \frac{1}{2}}{\displaystyle _{\mathrm{\Omega }_\epsilon }}K_1(\epsilon x)V_P^{}^2{\displaystyle \frac{1}{4}}{\displaystyle _{\mathrm{\Omega }_\epsilon }}K_2(\epsilon x)V_P^{}^4`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{(\mathrm{\Omega }Q^{})/\epsilon B_{\epsilon ^{1/4}}}{}}\left[|V^Q|^2+K_1(\epsilon x+Q^{})\left(V^Q\right)^2\right]`$
$`{\displaystyle \frac{1}{4}}{\displaystyle \underset{(\mathrm{\Omega }Q^{})/\epsilon B_{\epsilon ^{1/4}}}{}}K_2(\epsilon x+Q^{})\left(V^Q\right)^4+o(\epsilon )`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _^3}|V^Q|^2+{\displaystyle \frac{1}{2}}{\displaystyle _^3}K_1(Q)\left(V^Q\right)^2{\displaystyle \frac{1}{4}}{\displaystyle _^3}K_2(Q)\left(V^Q\right)^4`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle _^3}\left(K_1(\epsilon x+Q+\sqrt{\epsilon }e_1)K_1(Q)\right)\left(V^Q\right)^2`$
$`{\displaystyle \frac{1}{4}}{\displaystyle _^3}\left(K_2(\epsilon x+Q+\sqrt{\epsilon }e_1)K_2(Q)\right)\left(V^Q\right)^4+o(\epsilon )`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _^3}K_2(Q)\left(V^Q\right)^4+o(\epsilon ^{1/4})`$
$`=`$ $`{\displaystyle \frac{1}{2}}K_1(Q)^{\frac{1}{2}}K_2(Q)^1{\displaystyle _^3}W^4+o(\epsilon ^{1/4}).`$
Therefore
(4.8)
$$f_\epsilon ^J(V_P^{})=\frac{1}{2}K_1(Q)^{\frac{1}{2}}K_2(Q)^1_^3W^4+o(\epsilon ^{1/4}).$$
Now (4.4) follows immediately by (4.6), (4.7) and (4.8). $`\mathrm{}`$
## 5. A multiplicity result and proofs of theorems
In this section we give the proofs of our theorems. First of all, let us prove Theorem 1.1 as an easy consequence of the following multiplicity result:
###### Theorem 5.1.
Let (J) and (K) hold and suppose $`\mathrm{\Gamma }`$ has a compact set $`X\mathrm{\Omega }_0`$ where $`\mathrm{\Gamma }`$ achieves a strict local minimum (resp. maximum), in the sense that there exist $`\delta >0`$ and a $`\delta `$-neighborhood $`X_\delta \mathrm{\Omega }_0`$ of $`X`$ such that
$$b:=inf\{\mathrm{\Gamma }(Q):QX_\delta \}>a:=\mathrm{\Gamma }_{|_X},(\mathrm{resp}.sup\{\mathrm{\Gamma }(Q):QX_\delta \}<\mathrm{\Gamma }_{|_X}).$$
Then there exists $`\overline{\epsilon }>0`$ such that ($`๐ซ_\epsilon `$) has at least $`\mathrm{cat}(X,X_\delta )`$ solutions that concentrate near points of $`X_\delta `$, provided $`\epsilon (0,\overline{\epsilon })`$. Here $`\mathrm{cat}(X,X_\delta )`$ denotes the Lusternik-Schnirelman category of $`X`$ with respect to $`X_\delta `$.
Proof First of all, we fix $`\mu >0`$ in such a way that $`|Q|<\mu `$ for all $`QX`$. We will apply the finite dimensional procedure with such $`\mu `$ fixed.
We will treat only the case of minima, being the other one similar. We set $`Y=\{QX_\delta :๐_\epsilon (Q)c_0(a+b)/2\}`$, being $`c_0`$ defined in (4.5). By (4.4) it follows that there exists $`\overline{\epsilon }>0`$ such that
(5.1)
$$XYX_\delta ,$$
provided $`\epsilon (0,\overline{\epsilon })`$. Moreover, if $`QX_\delta `$ then $`\mathrm{\Gamma }(Q)b`$ and hence
$$๐_\epsilon (Q)c_0\mathrm{\Gamma }(Q)+o(\epsilon ^{1/4})c_0b+o(\epsilon ^{1/4}).$$
On the other side, if $`QY`$ then $`๐_\epsilon (Q)c_0(a+b)/2`$. Hence, for $`\epsilon `$ small, $`Y`$ cannot meet $`X_\delta `$ and this readily implies that $`Y`$ is compact. Then $`๐_\epsilon `$ possesses at least $`\mathrm{cat}(Y,X_\delta )`$ critical points in $`X_\delta `$. Using (5.1) and the properties of the category one gets
$$\mathrm{cat}(Y,Y)\mathrm{cat}(X,X_\delta ).$$
Moreover, by Lemma 4.1, we know that to critical points of $`๐_\epsilon `$ there correspond critical points of $`f_\epsilon `$ and so solutions of (2.1). Let $`Q_\epsilon X`$ be one of these critical points, if $`Q_\epsilon ^{}=Q_\epsilon +\sqrt{\epsilon }e_1`$, then
$$(u_\epsilon ^{Q_\epsilon },v_\epsilon ^{Q_\epsilon })=(U_{Q_\epsilon /\epsilon }+w(\epsilon ,Q_\epsilon ),V_{Q_\epsilon ^{}/\epsilon }+w^{}(\epsilon ,Q_\epsilon ))$$
is a solution of (2.1). Therefore
$`u_\epsilon ^{Q_\epsilon }(x/\epsilon )U_{Q_\epsilon /\epsilon }(x/\epsilon )=U^{Q_\epsilon }\left({\displaystyle \frac{xQ_\epsilon }{\epsilon }}\right)`$
$`v_\epsilon ^{Q_\epsilon }(x/\epsilon )V_{Q_\epsilon ^{}/\epsilon }(x/\epsilon )=V^{Q_\epsilon }\left({\displaystyle \frac{xQ_\epsilon ^{}}{\epsilon }}\right)`$
is a solution of ($`๐ซ_\epsilon `$) and also the concentration result follows. $`\mathrm{}`$
Let us now give a short proof of Theorem 1.5.
Proof of Theorem 1.5 We need only to observe that, in this case, the solutions of ($`\overline{๐ซ}_\epsilon `$) will be found near $`(\overline{U}^Q,\overline{V}^Q)`$, properly truncated, where $`\overline{U}^Q`$ is the unique solution of
$$\{\begin{array}{cc}\mathrm{\Delta }u+J_1(Q)u=J_2(Q)u^{2p1}\hfill & \text{in }^N,\hfill \\ u>0\hfill & \text{in }^N,\hfill \\ u(0)=\mathrm{max}_^Nu,\hfill & \end{array}$$
and $`\overline{V}_Q`$ is the unique solution of
$$\{\begin{array}{cc}\mathrm{\Delta }v+K_1(Q)v=K_2(Q)v^{2p1}\hfill & \text{in }^N,\hfill \\ v>0\hfill & \text{in }^N,\hfill \\ v(0)=\mathrm{max}_^Nv,\hfill & \end{array}$$
for an appropriate choice of $`Q\overline{\mathrm{\Omega }}_0`$. It is easy to see that
$`\overline{U}^Q(x)`$ $`=\left(J_1(Q)/J_2(Q)\right)^{1/(2p2)}\overline{W}\left(\sqrt{J_1(Q)}x\right),`$
$`\overline{V}^Q(x)`$ $`=\left(K_1(Q)/K_2(Q)\right)^{1/(2p2)}\overline{W}\left(\sqrt{K_1(Q)}x\right),`$
where $`\overline{W}`$ is the unique solution of
$$\{\begin{array}{ccc}\mathrm{\Delta }z+z=z^{2p1}\hfill & \text{in }^N,\hfill & \\ z>0\hfill & \text{in }^N,\hfill & \\ z(0)=\mathrm{max}_^Nz.\hfill & & \end{array}$$
At this point, we can repeat the previous arguments, with suitable modifications. $`\mathrm{}`$
###### Remark 5.2.
Of course, the analogous of Theorem 5.1 holds also for problem ($`\overline{๐ซ}_\epsilon `$).
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# On curvature and feedback classification of two-dimensional optimal control systems
## 1 Introduction
In Riemannian geometry the Gaussian curvature of a manifold reflects intrinsic properties of the geodesic flow, i.e. properties that do not depend on the choice of local coordinates. For example, the geodesics of the surface have no conjugate points if the curvature is non-positive. Indeed, these geodesics are extremals of a particular time optimal control problem the dynamics of which is given by
$$\dot{q}=\mathrm{cos}u๐_1(q)+\mathrm{sin}u๐_2(q),uS^1,$$
where $`(๐_1,๐_2)`$ forms an orthonormal frame of the Riemannian structure on the manifold. Our goal is to generalize the classical notion of Gaussian curvature of two-dimensional Riemannian surfaces for two-dimensional smooth optimal control problems. The notion of curvature tensor for non linear optimal control problems was first introduced in by A. A. Agrachev and R. V. Gamkrelidze with a purely variational description by means of Jacobi curves, which are curves in the Lagrangian Grassmannian. Here we will not deal with Jacobi curves but use the Cartanโs moving frame method in order to construct a feedback invariant frame associated to our optimal control problem and provide a less general but also more geometric definition of the curvature function.
Consider a control system of the form
$$\dot{q}=๐(q,u),qM,uU,$$
(1.1)
where $`M`$ and $`U`$ are smooth connected manifolds. Let $`\stackrel{\mathbf{~}}{๐}(\stackrel{~}{q},\stackrel{~}{u})`$, $`(\stackrel{~}{q},\stackrel{~}{u})\stackrel{~}{M}\times \stackrel{~}{U}`$, be the right-hand side of another such system. We say that the two systems are feedback-equivalent if there exists a diffeomorphism $`\mathrm{\Theta }:M\times U\stackrel{~}{M}\times \stackrel{~}{U}`$ of the form
$$\mathrm{\Theta }(q,u)=(\varphi (q),\psi (q,u))$$
(1.2)
which transforms the first system to the second, i.e. such that
$$T_q\varphi (๐(q,u))=\stackrel{\mathbf{~}}{๐}(\varphi (q),\psi (q,u)).$$
In the above diffeomorphism $`\varphi `$ plays the role of a change of coordinates in the state space $`M`$, and $`\psi `$ called pure feedback transformation reparametrizes the set of controls $`U`$ in a way depending on the state variable $`qM`$. Our aim is to provide feedback invariants for control system (1.1) when the manifold $`M`$ is of dimension two and the control set $`U`$ of dimension one what we suppose from now.
In this case, if the coordinates on the manifold are fixed, a control system of type (1.1) is parametrized by two functions of three variables, and the group of feedback transformations of type (1.2) is parametrized by two functions of two variables and one function of three variables. Therefore, we can a priori normalize only one function among the two functions defining control system (1.1). Thus, we expect to have only $`21=1`$ โprincipalโ feedback invariant, i.e. a function of three variables, in this equivalence problem.
All results of the present paper will be presented without proof. Anyway, most of these proofs can be found in the references cited at the end.
## 2 Curvature
Suppose that we want to minimize an integral cost $`_{t_0}^{t_1}\phi (q,u)๐t`$, along the trajectories of control system (1.1). We write the normal maximized Hamiltonian function of PMP (Pontryagin Maximum Principle) which is defined by
$$h(\lambda )=\underset{uU}{\mathrm{max}}\left(\lambda ,๐(q,u)\phi (q,u)\right),\lambda T_q^{}M,qM,$$
(2.1)
where $`,`$ denotes the canonical pairing between the tangent and the cotangent bundles over $`M`$. Hamiltonian $`h`$ is a function on the cotangent bundle $`T^{}M`$ which is, because of being independent of $`u`$, feedback-invariant. Thus, all objects construct from Hamiltonian $`h`$ through intrinsic relations will also be feedback invariants. As usual, if $`\sigma `$ denotes the symplectic two-form of the cotangent bundle over $`M`$, we define, via the relation $`\sigma (,\stackrel{\mathbf{}}{๐})=dh`$, the Hamiltonian vector field $`\stackrel{\mathbf{}}{๐}`$ associated to the Hamiltonian function $`h`$. Assume that $`h`$ is a smooth function, then the corresponding Hamiltonian vector field $`\stackrel{\mathbf{}}{๐}`$ is well-defined and tangent to the level set of $`h`$. PMP asserts (see e.g. ) that optimal trajectories of system (1.1) are projections onto $`M`$ of trajectories of the Hamiltonian system $`\dot{\lambda }=\stackrel{\mathbf{}}{๐}(\lambda )`$, in other words trajectories of Hamiltonian field $`\stackrel{\mathbf{}}{๐}`$ are extremals of our optimal control problem. Now fix a level set $`=h^1(e)`$ of our Hamiltonian, then the intersection $`_q=T_q^{}M`$ is a curve in the plane $`T_q^{}M`$ and under the regularity assumptions
$$๐(q,u)\frac{๐(q,u)}{u}0,\frac{๐(q,u)}{u}\frac{{}_{}{}^{2}๐(q,u)}{u^2}>0,qM,uU,$$
(2.2)
such a curve is a strictly convex curve surrounding the origin and it admits, up to sign and translation, a natural parameter providing us with a vector field $`๐_q`$ on $`_q`$ and by consequence with a vertical vector field $`๐`$ on $``$. Vector field $`๐`$ is characterized by the fact that it is, up to sign, the unique vector field on $``$ such that
$$L_๐^2s=s+bL_๐s,$$
(2.3)
where $`s`$ denotes the restriction to $``$ of Liouville one-form โ$`pdq`$โ of $`T^{}M`$ and $`b`$ is a smooth function on the level $``$.
Actually function $`b`$ is the feedback invariant of our control system that characterizes Riemannian problems. Namely control problem (1.1) defines a Riemannian geodesic problem if and only if the invariant $`b`$ is identically equal to zero.
Since vector fields $`\stackrel{\mathbf{}}{๐}`$ and $`๐`$ are feedback-invariant, it is natural to think that the curvature of our system may arise from a commutator relation of these fields. Indeed, the following theorem confirms this intuition.
###### Theorem 2.1.
Vector fields $`๐ฏ`$ and $`\stackrel{\mathbf{}}{๐ก}`$ satisfy the following nontrivial commutator relation:
$$[\stackrel{\mathbf{}}{๐},[๐,\stackrel{\mathbf{}}{๐}]]=\kappa ๐.$$
(2.4)
Proof of this theorem can be found in , . The coefficient $`\kappa `$ in the identity (2.4) is defined to be the curvature of our optimal control problem and since the fields $`\stackrel{\mathbf{}}{๐}`$ and $`๐`$ are feedback-invariant, the curvature $`\kappa `$ is also feedback-invariant.
Fix a system of local coordinates $`\lambda =(\theta ,q)`$ where $`\theta `$ parametrizes the fiber $`_q`$ so that $`๐=\frac{}{\theta }`$ and denote $`\frac{}{\theta }=^{}`$. Define the function $`c=c(\theta ,q)`$ by
$$d_qs=css^{},$$
(2.5)
where $`d_qs`$ is the differential of the Liouville one-form $`s`$ with respect to the horizontal coordinates. Then, the Hamiltonian field takes the form
$$\stackrel{\mathbf{}}{๐}=๐c\frac{}{\theta },$$
and the curvature $`\kappa `$ is evaluated as follows:
$$\kappa (\theta ,q)=L_\stackrel{\mathbf{}}{๐}^{}cL_\stackrel{\mathbf{}}{๐}c^{}.$$
(2.6)
###### Example 2.2.
Consider the control system corresponding to the geodesic problem on a two-dimensional Riemannian manifold:
$$\dot{q}=\mathrm{cos}u๐_1(q)+\mathrm{sin}u๐_2(q),uS^1.$$
(2.7)
In this case, control curvature $`\kappa `$ is the Gaussian curvature of the Riemannian manifold $`M`$ and it is evaluated as follows:
$$\kappa (q)=c_1^2c_2^2+L_{๐_1}c_2L_{๐_2}c_1,$$
(2.8)
where $`c_1`$, $`c_2`$ are the structural constants of the orthonormal frame $`(๐_1,๐_2)`$ on $`M`$:
$$[๐_1,๐_2]=c_1๐_1+c_2๐_2,c_1,c_2C^{\mathrm{}}(M).$$
See for the proof of this formula. Of course, for the Riemannian problem the curvature $`\kappa =\kappa (q)`$ depends only on the base point $`qM`$ as one can see from formula (2.8) but in general this is not the case: the curvature $`\kappa `$ depends also on the coordinate in the fiber $`_q`$ and thus is a function on the whole three-dimensional manifold $``$.
Observe that relations (2.3) and (2.4) define two feedback invariants: the function $`b`$ and the curvature $`\kappa `$. Both $`b`$ and $`\kappa `$ are functions on the three-dimensional level surface $``$, so that they are principal feedback invariants of our control system. Since our feedback equivalence problem admits only one invariant these functions are not โindependentโ. Indeed invariants $`b`$ and $`\kappa `$ are connected by the following differential relation:
$$L_๐\kappa +b\kappa +L_\stackrel{\mathbf{}}{๐}^2b=0,$$
(2.9)
which can easily be derived calculating some bracket relations between vector fields $`๐`$ and $`\stackrel{\mathbf{}}{๐}`$. In particular, relation (2.9) shows that in the special case of Riemannian problems, the curvature $`\kappa `$ is a function on the base manifold $`M`$ without any computation. Indeed, since Riemannian problems are characterized by the vanishing of function $`b`$, (2.9) reduces to $`L_๐\kappa =0`$.
## 3 Jacobi equation
It is easy to see that the regularity assumptions (2.2) imply $`\stackrel{\mathbf{}}{๐}๐[\stackrel{\mathbf{}}{๐},๐]0`$ so that vector fields $`\stackrel{\mathbf{}}{๐}`$, $`๐`$, $`[\stackrel{\mathbf{}}{๐},๐]`$ form a moving frame on the level surface $``$. In this section we use this moving frame to derive an ODE on conjugate time of our two-dimensional optimal control problem. This ODE, Jacobi equation in the moving frame, will show that the control curvature analogue to the Gaussian curvature enjoys similar properties.
Fix a point $`q_0M`$ and define a two-dimensional surface in $``$ by:
$$_0^t=e^{t\stackrel{\mathbf{}}{๐}}(_{q_0}),t,$$
where $`e^{t\stackrel{\mathbf{}}{๐}}`$ denotes the flow of the Hamiltonian field $`\stackrel{\mathbf{}}{๐}`$. The surface $`_0^t`$ is the lift in the cotangent bundle of trajectories $`tq(t)`$ in $`M`$ of control system (1.1) with starting point $`q(0)=q_0`$. We say that a point $`q=q(t)`$, $`t0`$, is conjugate to $`q_0`$ (or time $`t`$ is conjugate to zero) if $`q`$ is a critical value of the canonical projection
$$\pi :_0^tM.$$
(3.1)
It is easy to check that the tangent space $`T_\lambda _0^t`$, $`\lambda _0^t`$, is spanned by the vectors $`\stackrel{\mathbf{}}{๐}(\lambda )`$ and $`(e_{}^{t\stackrel{\mathbf{}}{๐}}๐)(\lambda )`$ so that the point $`q(t)=\pi (\lambda )`$ is conjugate to $`q_0`$ if and only if
$$(e_{}^{t\stackrel{\mathbf{}}{๐}}๐)(\lambda )\mathrm{span}(\stackrel{\mathbf{}}{๐}(\lambda ),๐(\lambda )).$$
Consider the decomposition of the vector field $`e_{}^{t\stackrel{\mathbf{}}{๐}}๐`$ in our moving frame on $``$:
$$e_{}^{t\stackrel{\mathbf{}}{๐}}๐=\alpha (t)\stackrel{\mathbf{}}{๐}+\beta (t)๐+\gamma (t)[\stackrel{\mathbf{}}{๐},๐].$$
It turns out that coefficients $`\alpha (t)`$, $`\beta (t)`$, $`\gamma (t)`$ are solutions to the Cauchy problem
$$\left(\begin{array}{c}\dot{\alpha }\\ \dot{\beta }\\ \dot{\gamma }\end{array}\right)=\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& \kappa _t\\ 0& 1& 0\end{array}\right)\left(\begin{array}{c}\alpha \\ \beta \\ \gamma \end{array}\right),\alpha (0)=1,\beta (0)=\gamma (0)=0,$$
(3.2)
where $`\kappa _t=\kappa (e^{t\stackrel{\mathbf{}}{๐}}(\lambda _0))`$, $`\pi (\lambda _0)=q_0`$. It is quite obvious that Cauchy problem (3.2) is equivalent to the second order linear ODE, called Jacobi equation
$$\ddot{\gamma }+\kappa _t\gamma =0,\gamma (0)=\gamma (t)=0,$$
(3.3)
where $`\kappa _t=\kappa (e^{t\stackrel{\mathbf{}}{๐}}(\lambda _0))`$. Thus an instant $`t`$ is a conjugate time for our optimal control problem if and only if there exists a non trivial solution to the boundary value problem (3.3). Using the Sturmโs comparison theorem for second order ODEs one can prove the following theorem about the occurrence of conjugate points for system (1.1).
###### Theorem 3.1.
Let $`q(t)`$, $`q(0)=q_0`$, be a solution of an optimal two-dimensional control problem and let $`\kappa _t`$ be the value of the curvature along an extremal $`\lambda (t)`$, $`\pi (\lambda (t))=q(t)`$.
1. If $`\kappa _t0`$ for all $`t0`$, then $`q_0`$ has no conjugate points for $`t[0,+\mathrm{}]`$.
2. If $`\kappa _t\kappa _1`$ (resp. $`\kappa _t<\kappa _1`$) for all $`t0`$, and some constant $`\kappa _1>0`$, then $`q_0`$ has no conjugate points along $`q()`$ for $`t[0,\pi /\sqrt{\kappa _1}[`$ (resp. for $`t[0,\pi /\sqrt{\kappa _1}]`$).
3. If $`0<\kappa _0\kappa _t`$ (resp. $`0<\kappa _0<\kappa _t`$), for all $`t0`$, then $`q_0`$ must have at least a conjugate point for $`t]0,\pi /\sqrt{\kappa _0}]`$ (resp. for $`t]0,\pi /\sqrt{\kappa _0}[`$).
The following theorem gives sufficient condition for a trajectory $`q(t)`$ on $`M`$ to be strongly locally optimal in terms of conjugate points (see for the definition of strong optimality and the proof of the following theorem).
###### Theorem 3.2.
Let the trajectory $`q(t)`$ be as in theorem 3.1. If the time interval $`]0,t_1]`$ does not contain conjugate points, then the trajectory $`q(t)`$ is strongly locally optimal for $`t[0,t_1]`$.
On the other hand if an instant $`t_c]0,t_1]`$ is conjugate to zero, then there exists an instant $`\stackrel{~}{t}]0,t_1]`$ where the trajectory $`q(t)`$, $`t]0,t_1]`$, ceases to be locally optimal.
###### Example 3.3.
Zermelo navigation problem (see for a detailed description). This problem is a time optimal control problem which consists of finding the quickest nautical path of a yacht in the presence of stationary sea currents. The sea surface is modeled by a two-dimensional Riemannian surface $`M`$ and the currents by an autonomous vector field $`๐ฟ\mathrm{Vec}M`$. Dynamics of optimal trajectories for Zermelo problem are given by
$$\dot{q}=๐ฟ(q)+\mathrm{cos}u๐_1+\mathrm{sin}u๐_2,qM,uS^1,$$
where $`(๐_1,๐_2)`$ form an orthonormal frame of the Riemannian surface $`M`$. Suppose that the manifold $`M`$ is the Euclidean plane $`^2`$. Then, in the coordinate system $`(q_1,q_2,u)`$ vector fields $`\stackrel{\mathbf{}}{๐}`$ and $`๐`$ read
$`\stackrel{\mathbf{}}{๐}`$ $`=`$ $`(๐ฟ_1+\mathrm{cos}u){\displaystyle \frac{}{q_1}}+(๐ฟ_2+\mathrm{sin}u){\displaystyle \frac{}{q_2}}D_q๐ฟ\left(\genfrac{}{}{0pt}{}{\mathrm{sin}u}{\mathrm{cos}u}\right),\left(\genfrac{}{}{0pt}{}{\mathrm{cos}u}{\mathrm{sin}u}\right){\displaystyle \frac{}{u}},`$
$`๐`$ $`=`$ $`\sqrt{๐ฟ_1(q)\mathrm{cos}u+๐ฟ_2(q)\mathrm{sin}u+1}{\displaystyle \frac{}{u}},`$
so that one can compute the curvature using formula (2.4) (here, $`,`$ denotes the scalar product between vectors). If we suppose moreover that the drift term is the linear field $`๐ฟ(q)=\left(\genfrac{}{}{0pt}{}{ab}{ba}\right)q`$ then, the control curvature for this problem is $`\kappa =a^2/4`$ (see ) and theorem 3.1 thus implies that there is no conjugate point along trajectories. The following theorem can also be proved using the definition of conjugate points.
###### Theorem 3.4.
There is no conjugate point for Zermelo navigation problem on $`^2`$ when the drift term $`๐`$ is a linear vector field.
Let us sketch the proof of this result (see for the detailed proof). Since the drift term is linear, it is easy to compute the map $`\pi `$ (see (3.1) for the definition) which takes the form:
$$\pi (t,u_0)=q(t,q_0,u_0)=e^{tA}q_0+_0^te^{(t\tau )A}\left(\genfrac{}{}{0pt}{}{\mathrm{cos}u(\tau )}{\mathrm{sin}u(\tau )}\right)๐\tau ,$$
where $`A`$ is the matrix representation of the linear drift term. Now, because the Hamiltonian flow preserves the Liouville one-form, it is easy to see that the differential $`d_{(t,u_0)}\pi `$ is of maximal rank if and only if $`d_{u_0}\pi `$ is of maximal rank. Saying this, it is now an easy task to find a vector $`v`$ such that $`d_{u_0}\pi (v)0`$ which completes the proof. The above theorem is also valid on $`^n`$ where the Zermelo navigation problem can be generalized without any difficulty; the proof is also similar.
## 4 Flat systems
In Riemannian geometry it is well-known that if the Gaussian curvature of the surface is nonzero then, one can not rectify simultaneously the geodesics by a change of coordinates. Only Riemannian flat systems, i.e. systems for which the geodesics are โstraight linesโ have this property. For control systems the situation is quite different first of all because control systems with zero curvature are not necessarily flat. We present here a new theorem which gives a characterization of flat control systems in terms of the feedback invariants $`\kappa `$ and $`b`$. We begin with the following definition.
###### Definition 4.1.
A control system $`\dot{q}=๐(q,u)`$ is said to be flat if it is feedback equivalent to a control system of the form $`\dot{q}=๐(u)`$.
It is obvious that a flat system has zero curvature but the contrary is in general not true. For example a Zermelo problem defined on the Euclidean plane $`^2`$ with a nonzero linear drift term is never flat.
Suppose that a control system satisfies
$$L_\stackrel{\mathbf{}}{๐}b=0.$$
(4.1)
The above property implies in particular that the plane curves $`_qT^{}M`$ are all of the same centro-affine length. Control systems of this type are very peculiar and have nice geometric properties that we do not discuss here. However such systems with zero curvature are characterized in the theorem below.
###### Theorem 4.2.
There exists a feedback transformation such that:
$$[๐(,u),\frac{๐(,u)}{u}]=0$$
(4.2)
if and only if the feedback invariants $`\kappa `$ and $`L_\stackrel{\mathbf{}}{๐ก}b`$ are identically equal to zero. Moreover if we fix local coordinates $`q=(q_1,q_2)`$ in $`M`$, then these systems can be parametrized by a one-parameter family of diffeomorphisms generated by the vector field:
$$๐ฟ_u=(a_1(u)+q_2)\frac{}{q_1}+(a_2(u,q_2)q_1)\frac{}{q_2}.$$
(4.3)
In the above theorem if $`u`$ is a control parameter such that the fields $`๐`$ and $`\frac{๐}{u}`$ commute then, vector field $`๐ฟ_u`$ is the infinitesimal generator of a diffeomorphism $`P_u\mathrm{Diff}M`$ such that
$$P_u(๐(,u),\frac{๐(,u)}{u})=(\left(\begin{array}{c}1\\ 0\end{array}\right),\left(\begin{array}{c}0\\ 1\end{array}\right)).$$
(4.4)
Notice that commutativity between vector fields $`๐`$ and $`\frac{๐}{u}`$ is not a feedback-invariant property. When the curvature is identically zero the above theorem shows that the PDE (4.1) can be reduced to the nonautononous ODE
$$\frac{dq}{du}=๐ฟ_u(q).$$
The following theorem characterizes flat control systems.
###### Theorem 4.3.
A control system of type (1.1) is flat if and only if its feedback invariants $`\kappa `$, $`L_\stackrel{\mathbf{}}{๐ก}b`$ and $`L_{[๐ฏ,\stackrel{\mathbf{}}{๐ก}]}b`$ vanish identically.
We do not discuss in detail proofs of theorems 4.2 and 4.3 in this paper but we roughly explain the main ideas. The proofs are based on the following differential equation which can easily be derived from the differentiation of the structural equations of our feedback invariant moving frame on $``$:
$$c^{\prime \prime }+bc^{}+c=L_\stackrel{\mathbf{}}{๐}b,$$
(4.5)
where $`c`$ is the function defined in (2.5). It follows immediately from this equation that if a control system is such that (4.2) holds (respectively if a control system is flat) then, its feedback invariants $`\kappa `$ and $`L_\stackrel{\mathbf{}}{๐}b`$ (respectively $`\kappa `$, $`L_\stackrel{\mathbf{}}{๐}b`$ and $`L_{[๐,\stackrel{\mathbf{}}{๐}]}b`$) vanish identically. To prove the converse observe first that if a control system has zero curvature then, the vector fields $`\stackrel{\mathbf{}}{๐}`$ and $`[๐,\stackrel{\mathbf{}}{๐}]`$ commute so that the choice of a natural parameter $`\theta `$ on the fibers $`_q`$ defines a foliation of $``$, the leaves of which are formed by the trajectories of the fields $`\stackrel{\mathbf{}}{๐}`$ and $`[๐,\stackrel{\mathbf{}}{๐}]`$. Now, choose the parameter $`\theta `$ (recall that this natural parameter is fixed only up to transformation of the form $`\theta \pm \theta +\varphi (q)`$) so that $`c`$ becomes zero which is possible since $`c`$ satisfies equation (4.5) with $`L_\stackrel{\mathbf{}}{๐}b=0`$. This shows in particular that there exists a feedback transformation so that (4.2) holds and by the Frobenius theorem one gets the existence of a diffeomorphism $`P_u\mathrm{Diff}M`$ such that (4.4) holds. In order to get the expression (4.3) we use Moserโs argument for which the key idea is to determine the diffeomorphisms $`P_u`$ by representing them as the flow of a family of vector fields $`๐ฟ_u`$ on $`M`$. We thus suppose that
$$\frac{d}{dt}P_u=๐ฟ_uP_u,P_{u_0}=\mathrm{Id},$$
and the expression of $`๐ฟ_u`$ in coordinates follows from differentiation with respect to $`u`$ of (4.4). This complete the proof of theorem 4.2. In order to complete the proof of theorem 4.3 one has just to check that $`L_\stackrel{\mathbf{}}{๐}b=0`$ and $`L_{[๐,\stackrel{\mathbf{}}{๐}]}b=0`$ imply that $`b=b(u)`$ which, in addition with (4.2) and $`\kappa =0`$ easily implies that the system is flat.
We now conclude our discussion with the following example.
###### Example 4.4.
Consider Zermelo navigation problem as in example 3.3. One can prove that this problem is flat if and only if the Riemannian surface $`M`$ is flat and the drift term $`๐ฟ`$ is constant.
### Acknowledgments
I am grateful to Professor Andrei A. Agrachev for fruitful discussions.
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# Quantum Tunneling Detection of Two-photon and Two-electron Processes
## Abstract
We analyze the operation of a quantum tunneling detector coupled to a coherent conductor. We demonstrate that in a certain energy range the output of the detector is determined by two-photon processes, two-electron processes and the interference of the two. We show how the individual contributions of these processes can be resolved in experiments.
The quantum nature of electron transfer in coherent conductors is seldom explicitly manifested in averaged current-voltage curves. To reveal it one should measure current noise and/or higher-order correlations of current comprising Full Counting Statistics which arise from the transfer Blanter and Bรผttiker (2000). Such measurements not only reveal the discrete nature of the charges transferred, they also quantify quantum many-body effects in electron transport and may be used for the detection of pairwise entanglement of transferred particles Nazarov (2003); Lorenzo and Nazarov (2005); Beenakker and Kindermann (2004); Samuelsson et al. (2004). If the noise is measured at frequencies in the quantum range, $`\mathrm{}\omega k_\mathrm{B}T`$, the measurement amounts to the detection of photons produced by the current fluctuations. This aspect is important in view of attempts to transfer quantum information from electrons to photons and back Cerletti et al. (2004).
It was demonstrated in Gavish et al. (2000) that one needs a quantum detector to measure quantum noise. Indeed, any classical measurement of a fluctuating quantity would give a noise spectrum symmetric in frequency, $`S(\omega )=S(\omega )`$. A quantum tunneling detector is generally a quantum two-level system with a level separation $`\epsilon >0`$. The result of detection are two transition rates: $`\mathrm{\Gamma }_{\mathrm{up}}`$ from the lower to the higher level and $`\mathrm{\Gamma }_{\mathrm{down}}`$ for the reverse direction. The most probable transitions are accompanied by either absorption or emission of a photon of matching energy $`\mathrm{}\omega =\epsilon `$. One can define the noise spectrum in such a way that it is proportional to the transition rates $`S(\pm \epsilon /\mathrm{})\mathrm{\Gamma }_{\mathrm{up},\mathrm{down}}(\epsilon )`$. Differences between $`\mathrm{\Gamma }_{\mathrm{up},\mathrm{down}}`$ thus manifest the quantum nature of noise. If the source of noise is a coherent conductor biased by a voltage $`V`$, detector signals in the range $`\epsilon <eV`$ are readily interpreted in terms of single electron transfers through the conductor. The maximum energy gain available for electrons in the course of such transfer is $`eV`$. Consequently this value also limits the energy of the emitted photon.
A first proposal for the experimental realization of a quantum tunneling detector included transitions between two localized electron states in a double quantum dot Aguado and Kouwenhoven (2000). However it does not matter much if the tunneling occurs between localized or delocalized electron states and if all tunnel events are accompanied by the same energy transfer $`\epsilon `$. In most practical cases the energy dependence of the rates $`\mathrm{\Gamma }_{\mathrm{up},\mathrm{down}}`$ can be readily extracted from the measurement results. This is why quantum tunneling detection has been experimentally realized in a superconducting tunnel junction Deblock et al. (2003) and in a single quantum dot Balestro et al. .
In this Letter we study quantum tunneling detection in the range $`eV<\epsilon <2eV`$ assuming $`\epsilon ,eVk_\mathrm{B}T`$. The motivation is that for these energies the detector is not sensitive to single-electron one-photon processes described above and its output โ the transition rate $`\mathrm{\Gamma }_{\mathrm{up}}`$โ is determined by much more interesting two-particle processes. It is clear from plain energy considerations that transitions may originate from two-photon processes. Such two-photon absorption can occur given any non-equilibrium photon distribution bounded by $`eV`$, not necessarily produced by a coherent conductor. Less obvious and specific for a coherent conductor is a cooperative two-electron process. Indeed, if two electrons team up in crossing the conductor they can emit a single photon with an energy up to $`2eV`$. Essential for this cooperation are electron-electron interactions. It is known Ingold and Nazarov (1992) that the most important electron-electron interaction in this energy range is due to the electromagnetic environment of the conductor, the same environment in which the non-equilibrium photons dwell.
We quantify the signals due to two-photon and two-electron events and find them to be of the same order of magnitude. We also show that part of the signal is due to quantum interference of these two processes: one-and-half-photon absorption events. We demonstrate how different contributions can be separated in experiments thereby facilitating the direct observation of two-particle processes in the context of quantum transport.
We concentrate on a model circuit consisting of four elements as given in figure 1. A voltage biased coherent contact characterized by a set of transmission eigenvalues $`\{T_n\}`$ is embedded in an electromagnetic environment with impedance $`Z_\omega `$. The environment transforms the current fluctuations in the conductor to voltage fluctuations in node A which are conveniently expressed in terms of a phase $`\phi =\frac{e}{\mathrm{}}๐tV(t)`$. The most general model including the detector and coherent contact would be a four-pole circuit studied in Kindermann et al. (2004) that couples two poles of the detector with two poles of the contact. Here, we concentrate on the experimentally relevant case of capacitive coupling. Owing to voltage division between two capacitors, the detector senses a fraction $`0<\alpha <1`$ of the voltage fluctuations in node A. We will see that changing the โvisibilityโ parameter $`\alpha `$, enables the separation of two-electron and two-photon processes in experiments. The relevant impedance is made up of an environmental impedance combined with that of two capacitors and that of the coherent contact. We measure this impedance $`z_\omega `$ in units of $`R_K2\pi \mathrm{}/e^2`$ and assume the low-impedance limit $`z_\omega 1`$; this provides us with a physically justified small parameter.
The detector consists of two localized charge states connected by a tunnel amplitude $`๐ฏ`$. In the presence of voltage fluctuations in the node A, the amplitude is modified as follows: $`๐ฏ(t)๐ฏe^{i\alpha \phi (t)}`$. In perturbation theory, the inelastic tunneling rate between two states separated by $`\epsilon `$ is given by correlators of $`\alpha \phi (t)`$ Ingold and Nazarov (1992)
$$\mathrm{\Gamma }(\epsilon )=\frac{|๐ฏ|^2}{2\pi \mathrm{}^2}๐te^{i\alpha \phi (t)}e^{i\alpha \phi (0)}e^{\frac{i}{\mathrm{}}\epsilon t}.$$
(1)
The rate $`\mathrm{\Gamma }(\epsilon )`$ is therefore the Fourier transform of the correlation function $`e^{i\alpha \phi (t)}e^{i\alpha \phi (0)}`$, $`\alpha \phi (t)`$ being the phase fluctuations over the detector. From now on we take $`\mathrm{}=e=k_\mathrm{B}=1`$.
Equation (1) tells us that the inelastic tunneling rates in the detector are completely determined by the voltage fluctuations over the junction. Therefore measuring the inelastic current through the dots we are sensitive to the noise spectrum of the environment.
To evaluate $`e^{i\alpha \phi (t)}e^{i\alpha \phi (0)}`$ we construct a path integral representation of this quantity using a non-equilibrium Keldysh technique Rammer and Smith (1986) for quantum-circuits Kindermann et al. (2003)
$$\begin{array}{cc}\hfill e^{i\alpha \phi (t)}e& {}_{}{}^{i\alpha \phi (0)}=๐[\mathit{\varphi }]\mathrm{exp}\{iS_{\mathrm{env}}[\mathit{\varphi }]\hfill \\ & iS_{\mathrm{cond}}[\mathit{\varphi }]+i\alpha [\phi ^+(0)+\phi ^{}(t)]\}.\hfill \end{array}$$
(2)
The integration goes over the time-dependent fluctuating fields $`\phi ^\pm (t)`$ in node $`A`$, $`\pm `$ corresponding to the forward (backward) part of the Keldysh contour. $`S_{\mathrm{env}}`$ and $`S_{\mathrm{cond}}`$ are the contributions to the Keldysh action originating from the environment and the coherent conductor respectively.
Since the environment is linear, its action is quadratic in the fields and at zero temperature reads (cf. Kindermann and Nazarov (2003))
$$S_{\mathrm{env}}=๐\omega \mathit{\varphi }_\omega ^TA(\omega )\mathit{\varphi }_\omega $$
(3)
with
$$A(\omega )=\frac{i}{2}\left(\begin{array}{cccccccccccccccccccc}0& \frac{\omega }{z_\omega }& & & & & & & & & & & & & & & & & & \\ \frac{\omega }{z_\omega }& |\omega |\mathrm{Re}\{\frac{1}{z_\omega }\}& & & & & & & & & & & & & & & & & & \end{array}\right),$$
$`z_\omega `$ being the corresponding impedance. We use Fourier transformed fields $`\mathit{\varphi }_\omega =(\varphi _\omega ,\chi _\omega )^T`$ defined with $`\phi ^\pm =\varphi \pm \frac{1}{2}\chi `$.
All non-quadratic contributions to the action originate from the coherent conductor. The action $`S_{\mathrm{cond}}`$ can be expressed in terms of Keldysh Green functions $`\stackrel{ห}{G}_{L,R}`$ of electrons in the reservoirs left and right of the contact Kindermann and Nazarov (2003)
$$S_{\mathrm{cond}}=\frac{i}{2}\underset{n}{}\mathrm{Tr}\mathrm{ln}[1+\frac{T_n}{4}(\{\stackrel{ห}{G}_L(\mathit{\varphi }),\stackrel{ห}{G}_R\}2)].$$
(4)
The fields $`\mathit{\varphi }`$ enter in this action via the gauge transform of $`\stackrel{ห}{G}_L`$.Kindermann and Nazarov (2003)
To comprehend the physics involved, let us first disregard any non-quadratic parts and take only the quadratic part of $`S_{\mathrm{cond}}`$. In this case the path integral is Gaussian, and can be evaluated exactly. We recover the well known result from $`P(E)`$-theory (cf. Aguado and Kouwenhoven (2000); Ingold and Nazarov (1992)): $`e^{i\alpha \phi (t)}e^{i\alpha \phi (0)}=\mathrm{exp}[J(t)]`$ with
$$J(t)=\alpha \phi (t)\alpha \phi (0)=\alpha ^2๐\omega \frac{|z_\omega |^2}{\omega ^2}K(\omega )[e^{i\omega t}1].$$
(5)
The impedance includes the dimensionless conductance $`g_c_nT_n`$ of the contact.
In the limit of $`T=0`$ we find in agreement with results of Gavish et al. (2000); Aguado and Kouwenhoven (2000)
$$\begin{array}{cc}\hfill K(\omega )=g_c& \{FD(\omega +V)+(22F)D(\omega )\hfill \\ & +FD(\omega V)\}+2\mathrm{R}\mathrm{e}\{\frac{1}{z_\omega }\}D(\omega )\hfill \end{array}$$
(6)
with $`D(\omega )\omega \theta (\omega )`$ and the Fano factor $`F_nT_n(1T_n)/_nT_n`$. The first term in $`K(\omega )`$ ($`g_c`$) represents the non-equilibrium current noise spectrum of the coherent contact that vanishes for $`\omega >V`$. In physical terms this means that the highest energy $`\omega `$ an electron can emit traversing the conductor is exactly $`V`$. The second part represents the spectrum of the environment. It is zero for $`\omega >0`$, since the environment can only absorb energy at $`T=0`$.
The time-dependent part of $`J(t)`$ is the Fourier transform of $`K(\omega )|z_\omega |^2/\omega ^2`$ and $`\mathrm{\Gamma }(\epsilon )`$ is in turn the Fourier transform of $`\mathrm{exp}[J(t)]`$. If we expand $`\mathrm{exp}[J(t)]`$ in terms of $`J(t)`$, the $`n`$-th term presents the contribution of a process involving absorption of $`n`$ photons in the detector. Such an $`n`$-photon process dominates in the interval $`(n1)V<\epsilon <nV`$ and its contribution is proportional to $`\alpha ^{2n}`$.
The one-photon contribution gives $`\mathrm{\Gamma }_{\mathrm{up}}^{(1)}/\mathrm{\Gamma }_{\mathrm{down}}zg_cF`$. Each extra photon brings in a small factor, such that $`\mathrm{\Gamma }^{(n+1)}/\mathrm{\Gamma }^{(n)}\alpha ^2z^2g_cF`$. This is seen as a staircase in the log plot presented in Figure 2.
What we did was wrong since we did not take into account the non-quadratic terms in the action. These describe more interesting many-electron processes and areโas we show belowโof the same order of magnitude. Since the path integral in (2) can not be evaluated in general, we proceed by perturbative expansion.
Indeed, since $`|z|1`$, the Gaussian part of the action, being proportional to $`z^1`$, suppresses fluctuations in the path integral and we can treat the remaining part perturbatively. First we expand $`iS_{\mathrm{cond}}[\mathit{\varphi }]`$ around $`\mathit{\varphi }=0`$. As mentioned previously (see the discussion below equation (5)), the first and second order terms just renormalize the impedance. The exponential of the remaining higher order terms is then again expanded in $`\mathit{\varphi }`$ around $`\mathit{\varphi }=0`$. This expansion may be represented in terms of diagrams such as those in figure 3. Diagram (a) represents a high order term, from which the general structure becomes clear: Diagrams consist of lines, polygons and external vertices. The expansion contains not only connected diagrams, but all disconnected diagrams as well. A polygon with $`n`$ vertices is associated with the symmetrized $`n`$-th order coefficient in the Taylor expansion of $`iS_{\mathrm{cond}}[\mathit{\varphi }]`$. Each polygon contributes a factor $`g_c`$. Lines represent propagators of $`\mathit{\varphi }`$ corresponding to the Gaussian action with renormalized impedance. They are of order $`z`$ making $`n`$-line diagrams $`z^n`$ in leading order. External vertices (indicated by dots in the figure) are associated with the time-dependent linear term $`i\alpha [\phi ^+(0)+\phi ^{}(t)]`$ in equation (2). Thus a diagram with $`s`$ external vertices gives a correction proportional to $`\alpha ^s`$. Furthermore, diagrams without external vertices are time-independent and according to equation (1) contribute only to elastic tunneling processes. Diagrams (b) to (f) represent some of the lowest order terms in the expansion.
We consider transitions in the detector where energies between $`V`$ and $`2V`$ are absorbed. In this interval, the detector output is given by diagrams (b), (c) and (d), which are proportional to $`\alpha ^4,\alpha ^3`$ and $`\alpha ^2`$ respectively.
From the results presented in figure 2 we have learned that $`n`$-photon processes come with a coefficient $`\alpha ^{2n}`$. Hence the $`\alpha ^3`$ contribution is not readily expected: it seems to signal a process with one-and-half photons absorbed. We disregard diagram (e) which only contributes to elastic processes. In the energy interval considered, the combined $`z^3`$ contribution of the included diagrams is zero and we obtain a tunneling rate that goes as $`g_c^2z^4`$. Since a diagram like (f) has four lines, it could potentially contribute to the current with the same order in $`z`$. However, its contribution can only be proportional to $`g_c`$ and is disregarded.
The expansion of $`S_{\mathrm{cond}}`$ up to fourth order terms and subsequent evaluation of the diagrams is straightforward but requires rather involved and lengthy calculations. Fortunately in the interval of interest the three contributions can be combined in a compact expression
$`\mathrm{\Gamma }_{\mathrm{up}}=2|๐ฏ|^2g_c^2F^2{\displaystyle _{\epsilon V}^V}๐\omega (V\omega )(\epsilon V\omega )`$
$`{\displaystyle \frac{|z_\omega |^2}{\omega ^2}}\left|{\displaystyle \frac{\alpha ^2}{2}}{\displaystyle \frac{z_{\epsilon \omega }}{\epsilon \omega }}+\alpha {\displaystyle \frac{z_\epsilon }{\epsilon }}\right|^2,`$ (7)
which is the main result of our work. The rate is proportional to the square of the zero-frequency current noise $`S_{\mathrm{cond}}(0)=\frac{2}{\pi }g_cF`$.
The part proportional to $`\alpha ^4`$ (diagram b) represents a two-photon process originating from the quadratic part of $`S_{\mathrm{cond}}`$ and was already present in figure 2. We have thus shown that there are contributions of the same order resulting from non-linearities in the conductor. The $`\alpha ^2`$ term (diagram d) is the result of the two-electron and one-photon process expected from general reasoning presented in the introduction. We see that the $`\alpha ^3`$ term comes from the cross-term in the modulus square. This unambiguously identifies digram (c) as the result of quantum interference of the two-electron process and the two-photon processesโ an interpretation that was not obvious from the beginning.
To understand this interference, we note that the photon modes involved are delocalized across the whole circuit. A photon in each mode can be absorbed in the detector as well as in the environment or the contact. An elementary process is such that the final state differs from the initial state by two photons absorbed in two given modes. The final state can be reached by two amplitudes: one with both photons absorbed in the detector and one with a photon absorbed in the detector and a photon absorbed in the environment. While the squares of these amplitudes represent the probabilities of two-photon and two-electron processes respectively, their cross-term gives rise to an interference contribution $`\alpha ^3`$.
The simplest concrete model is that of a frequency-independent impedance, $`z_\omega =z`$ at $`\omega V`$. The integration in Eq. 7 yields for the three distinct contributions ( $`\stackrel{~}{\epsilon }=\epsilon /V`$, $`1<\stackrel{~}{\epsilon }<2`$ )
$$\frac{\mathrm{\Gamma }_i}{z^4|๐ฏ|^2g_c^2F^2}=\{\begin{array}{c}\alpha ^4\left[\frac{22\stackrel{~}{\epsilon }+\stackrel{~}{\epsilon }^2}{\stackrel{~}{\epsilon }^3}\mathrm{ln}(\stackrel{~}{\epsilon }1)\frac{2\stackrel{~}{\epsilon }}{\stackrel{~}{\epsilon }^2}\right]\hfill \\ 2\alpha ^3\left[\frac{22\stackrel{~}{\epsilon }+\stackrel{~}{\epsilon }^2}{\stackrel{~}{\epsilon }^3}\mathrm{ln}(\stackrel{~}{\epsilon }1)\frac{2\stackrel{~}{\epsilon }}{\stackrel{~}{\epsilon }^2}\right]\hfill \\ \alpha ^2\left[\frac{2}{\stackrel{~}{\epsilon }}\mathrm{ln}(\stackrel{~}{\epsilon }1)\frac{4(2\stackrel{~}{\epsilon })}{\stackrel{~}{\epsilon }^2}\right]\hfill \end{array}$$
(8)
All contributions scale as $`(\epsilon 2V)^3`$ at the two-photon threshold and logarithmically diverge at approaching the one-photon threshold. (see Fig. 4).
Here, we have quantified the contributions for a very specific non-linear quantum noise source: a coherent conductor. However in the case of any unknown source of this kind the $`\alpha `$-dependence of the contributions allows one to separate and identify them experimentally (right pane of Fig. 4). One would measure the detector output changing the coupling to the detector. Formally, three measurements at three different $`\alpha `$ are sufficient to determine the relative strength of all three contributions. In any case, in the limit of small coupling $`\alpha 0`$ the detector output is dominated by two-electron events. Further characterization may be achieved by engineering of a frequency-dependent impedance. For instance, setting $`z(\omega =ฯต)`$ to 0 kills both interference and two-electron contribution.
To conclude, we have shown that the quantum tunneling detector in the energy interval specified is selectively sensitive to two-particle processes. The detector output is generally determined by three contributions: two-photon processes, two-electron processes and the interference of the two. These three sources can be distinguished experimentally by measuring at different couplings $`\alpha `$ to the detector. Our results thus facilitate the direct observation of many-particle events in the context of quantum transport. This work was supported by the Dutch Foundation for Fundamental Research on Matter (FOM).
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# One-loop corrections to the Drell-Yan process in SANC
## 1 Introduction
Precision studies of the Drell-Yan process are vitally important for high energy hadronic colliders. This process provides information about weak interactions and contributes to the background to many of the searches for physics beyond the Standard Model. One-loop QED and electroweak (EW) radiative corrections (RC) to the Drell-Yan process at high energy hadronic collider were calculated by several groups in the past, see papers Mosolov:1981xk ; Soroko:1990ug ; Wackeroth:1996hz ; Baur:1998kt ; Dittmaier:2001ay ; Baur:2004ig and references therein. Here we present the results for the corrections to the charged current Drell-Yan process, obtained within the automatized system SANC Andonov:2004hi ; SANCwww and some comparisons with earlier calculations. Starting from the construction of helicity amplitudes and EW form factors, SANC performs calculation of the process cross section and produces computer codes, which can be further used in the experimental data analysis.
## 2 Preliminaries and Notation
Let us start with the partonic level, where we will consider interactions of free quarks (partons). The differential Born-level cross section of the process
$`\overline{d}(p_1)+u(p_2)l^+(p_4)+\nu _l(p_3)`$ (1)
in the center-of-mass system of the initial quarks reads
$`{\displaystyle \frac{\mathrm{d}\widehat{\sigma }_0}{\mathrm{d}\widehat{\mathrm{\Omega }}}}={\displaystyle \frac{1}{4}}{\displaystyle \frac{1}{N_c}}|V_{ud}|^2{\displaystyle \frac{G_F^2M_W^2}{2\pi \widehat{s}}}{\displaystyle \frac{\widehat{u}^2}{(\widehat{s}M_W^2)^2+\mathrm{\Gamma }_W^2(\widehat{s})M_W^2}},`$
$`\widehat{s}=(p_1+p_2)^2,\widehat{u}=(p_1p_3)^2,`$ (2)
where $`N_C=3`$ is the number of quark colors; $`V_{ud}`$ is the relevant element of the CKM matrix; $`G_F`$ is the Fermi coupling constant; $`M_W`$ and $`\mathrm{\Gamma }_W`$ are the mass and the width of the $`W`$-boson, respectively.
## 3 Radiative Corrections at the Partonic Level
In order to get a more accurate description of the process we should go beyond the Born approximation and take into account different sources of radiative corrections. Here we will consider only EW contributions to the corrections, while effects of higher order QCD contributions (and mixed effects) are left beyond the scope of our study.
As usually, we subdivide the EW RC into the virtual (loop) ones, the ones due to soft photon emission, and the ones due to hard photon emission. An auxiliary parameter $`\overline{\omega }`$ separates the soft and hard photonic contributions.
In the automatized system SANCwww , the virtual corrections are accessible via menu chain SANC $``$ EW $``$ Processes $``$ 4 legs $``$ 4f $``$ Charged Current $``$ f1 f1โ $``$ f fโ (FF). The module, loaded at the end of this chain computes on-line the scalar form factors of the partonic sub-process (1). The parallel module โฆf1 f1โ $``$ f fโ (HA) provides the relevant helicity amplitudes. For more details see Section 2.5 of the SANC description Andonov:2004hi and the book Bardin:1999ak .
The real photon emission process
$`\overline{d}(p_1)+u(p_2)l^+(p_4)+\nu _l(p_3)+\gamma (p_5)`$ (3)
should be taken into account as well. Integration over the phase space in this case can be performed either (semi-)analytically or by means of a Monte Carlo integrator.
The first possibility is realized within the SANC environment. Now we have there two branches. The first one contains the complete chain of analytical integrals over the hard photon phase space. It provides at the partonic level the double-differential distribution $`\mathrm{d}^2\widehat{\sigma }_{\mathrm{hard}}/(\mathrm{d}c\mathrm{d}\widehat{s}^{})`$ and the single differential distribution $`\mathrm{d}\widehat{\sigma }_{\mathrm{hard}}/\mathrm{d}c`$, where $`c=\mathrm{cos}\mathrm{}(\stackrel{}{p}_2\stackrel{}{p}_4)`$ and $`\widehat{s}^{}=(p_3+p_4)^2`$. The second branch provides the double-differential distribution $`\mathrm{d}^2\widehat{\sigma }_{\mathrm{hard}}/(\mathrm{d}c\mathrm{d}M_x^2)`$, where $`M_x^2=2p_3p_5`$ which is directly related to the charged lepton energy in the center-of-mass system of the initial quarks:
$`\widehat{E}_\mu =p_4^0={\displaystyle \frac{\widehat{s}+m_l^2M_x^2}{2\sqrt{\widehat{s}}}}.`$ (4)
We managed also to obtain analytically the hard photon contribution as the single differential distribution $`\mathrm{d}\widehat{\sigma }_{\mathrm{hard}}/\mathrm{d}\widehat{s}^{}`$. In this case we can use a system of reference with the $`z`$-axis along the real photon momentum $`\stackrel{}{p}_5`$. There the integration over three angular variables is rather easy and we have a possibility even to keep all the light masses exactly. Below we give the expression without mass terms because of its simplicity:
$`{\displaystyle \frac{\mathrm{d}\widehat{\sigma }_{\mathrm{hard}}}{\mathrm{d}\widehat{s}^{}}}=\widehat{\sigma }_0{\displaystyle \frac{\alpha }{2\pi }}{\displaystyle \frac{1}{\widehat{s}^2}}{\displaystyle \frac{1}{\widehat{s}\widehat{s}^{}}}\{[\widehat{s}^2+\widehat{s}^2][Q_l^2`$
$`\times \left(\mathrm{ln}{\displaystyle \frac{\widehat{s}^{}}{m_l^2}}1\right)+{\displaystyle \frac{\widehat{s}^{}}{\widehat{s}}}{\displaystyle \frac{(\widehat{s}M_W^2)^2+\mathrm{\Gamma }_W^2M_W^2}{(\widehat{s}^{}M_W^2)^2+\mathrm{\Gamma }_W^2M_W^2}}`$
$`\times [Q_u^2(\mathrm{ln}{\displaystyle \frac{\widehat{s}}{m_u^2}}1)+Q_d^2(\mathrm{ln}{\displaystyle \frac{\widehat{s}}{m_d^2}}1)]]`$
$`{\displaystyle \frac{2}{3}}(\widehat{s}^2+\widehat{s}\widehat{s}^{}+\widehat{s}^2)[Q_l^2`$
$`+{\displaystyle \frac{\widehat{s}^{}}{\widehat{s}}}{\displaystyle \frac{(\widehat{s}M_W^2)^2+\mathrm{\Gamma }_W^2M_W^2}{(\widehat{s}^{}M_W^2)^2+\mathrm{\Gamma }_W^2M_W^2}}]`$
$`{\displaystyle \frac{1}{3}}\left[\widehat{s}^{}\left(\widehat{s}+\widehat{s}^{}\right)\right]Q_l\left(4Q_u+5Q_d\right)`$
$`\times {\displaystyle \frac{\left(M_W^2\widehat{s}\right)\left(M_W^2\widehat{s}^{}\right)+\mathrm{\Gamma }_W^2M_W^2}{(\widehat{s}^{}M_W^2)^2+\mathrm{\Gamma }_W^2M_W^2}}\},`$ (5)
where $`Q_l`$, $`Q_u`$, and $`Q_d`$ are the charges of the charged lepton, up-quark and down-quark, respectively.
The differential distributions of the tree-level radiative process $`\overline{d}+u\mu ^++\nu _\mu +\gamma `$ were compared with the corresponding distributions obtained by means of the CompHEP package Boos:2004kh . Cross section distributions in the cosine of the outgoing muon angle and in the muon energy are considered. 20 bins are constructed for each of the distributions. Bins in the muon energy are
$`(n_{\mathrm{bin}}1)\times 5\mathrm{GeV}<E_\mu <n_{\mathrm{bin}}\times 5\mathrm{GeV}.`$ (6)
The cut on the muon energy ($`E_\mu <95`$ GeV) is imposed in both distributions to avoid the region with soft photons, where CompHEP is not supposed to work well. The angular bins are
$`1+{\displaystyle \frac{n_{\mathrm{bin}}1}{10}}<c<1+{\displaystyle \frac{n_{\mathrm{bin}}}{10}}.`$ (7)
The $`\alpha (M_Z)`$ electroweak scheme (realized according to the CompHEP conventions) was used. An agreement was found as can be seen from Table 1.
For the two choices of variables we have simple analytical expressions for the corresponding soft photon contributions. The infrared singularities in them are regularized by the auxiliary photon mass. The energy of a soft photon is limited from above by a cut in the integral either over $`\widehat{s}^{}`$ or over $`M_x^2`$.
In order to have the possibility to impose experimental cuts and event selection procedures of any kind, we can use a Monte Carlo integration routine based on the Vegas algorithm Lepage:1977sw . In this case we perform a 4(6)-fold numerical integration to get the hard photon contribution to the partonic (hadronic) cross section. To get the total EW correction we add also the contributions of the soft photon emission and the ones of the virtual loops. The cancellation of the dependence on the auxiliary parameter $`\overline{\omega }`$ in the sum is observed numerically.
Using the splitting of the $`W`$-boson propagators in the case of real photon emission off the virtual $`W`$, we separate the contributions of the initial state radiation, the final state one, and their interference in a gauge invariant way Berends:1984qa . The splitting is introduced by the following formula:
$`{\displaystyle \frac{1}{\widehat{s}(M_W\mathrm{i}\mathrm{\Gamma }_W)^2}}{\displaystyle \frac{1}{\widehat{s}^{}(M_W\mathrm{i}\mathrm{\Gamma }_W)^2}}`$
$`={\displaystyle \frac{1}{(\widehat{s}\widehat{s}^{})}}({\displaystyle \frac{1}{\widehat{s}^{}(M_W\mathrm{i}\mathrm{\Gamma }_W)^2}}`$
$`{\displaystyle \frac{1}{\widehat{s}(M_W\mathrm{i}\mathrm{\Gamma }_W)^2}}).`$ (8)
In the center-of-mass system $`(\widehat{s}\widehat{s}^{})=2p_5^0\sqrt{\widehat{s}}`$. The fixed $`W`$-width scheme is used here and in what follows.
In the course of calculations of the $`๐ช\left(\alpha \right)`$ corrections we met the so-called on-shell singularities, which appear in the form of $`\mathrm{ln}(\widehat{s}M_W^2+iฯต)`$. As was shown in detail in Ref. Wackeroth:1996hz , they can be regularized by the $`W`$-width:
$`\mathrm{ln}(\widehat{s}^{}M_W^2+iฯต)\mathrm{ln}(\widehat{s}^{}M_W^2+iM_W\mathrm{\Gamma }_W).`$ (9)
In the analytical formulae for radiative corrections one can find logarithms with quark and lepton mass singularities:
$`\mathrm{ln}{\displaystyle \frac{\widehat{s}}{m_l^2}},\mathrm{ln}{\displaystyle \frac{\widehat{s}}{m_u^2}},\mathrm{ln}{\displaystyle \frac{\widehat{s}}{m_d^2}}.`$ (10)
In the experimental set-up with calorimetric registration of the final state charged particles (typical for electrons), the lepton mass singularity cancels out in the result for the correction to an observable cross section in accordance with the KinoshitaโLeeโNauenberg theorem Kinoshita:1962ur ; Lee:1964is . But if the experiment is measuring the energy of the charged lepton without summing it with the energies of accompanying collinear photons (typical for muons), the logarithms with the lepton mass singularity remain in the result and give a considerable numerical contribution. Re-summation of these logs in higher orders was discussed in Refs. Placzek:2003zg ; CarloniCalame:2003ux ; CarloniCalame:2003ck .
### 3.1 Treatment of Quark Mass Singularities
One-loop radiative corrections contain terms proportional to the logarithms of the quark masses, $`\mathrm{ln}(\widehat{s}/m_{u,d}^2)`$. They come from the initial state radiation contributions including hard, soft and virtual photon emission. Such initial state mass singularities are well known, for instance, in the process of $`e^+e^{}`$ annihilation. But in the case of hadron collisions these logs have been already effectively taken into account in the parton density functions (PDFโs). In fact, in the procedure of PDFโs extraction from the experimental data, QED radiative corrections to the quark line have not been systematically subtracted. Therefore the present PDFโs effectively include not only the QCD evolution but also the QED one. Moreover, it is known that the leading log behaviors of the QED and QCD DGLAP evolution of quark density functions are similar (proportional to each other). So one gets evolution of PDFโs with an effective coupling constant
$`\alpha _s^{\mathrm{eff}}\alpha _s+{\displaystyle \frac{Q_i^2}{C_F}}\alpha ,`$ (11)
where $`\alpha _s`$ is the strong coupling constant, $`\alpha `$ is the fine structure constant, $`Q_i`$ is the quark charge, and $`C_F`$ is the QCD color factor. The nontrivial difference between the QED evolution and the QCD one starts to appear in higher orders, and the corresponding numerical effect is small compared to the remaining QCD uncertainties in PDFโs Kripfganz:bd ; Spiesberger:1994dm ; Roth:2004ti ; Martin:2004dh . The best approach to the whole problem would be to reโanalyze all the experimental DIS data taking into account QED corrections to the quark line at least at the nextโtoโleading order. But for the present moment we can limit ourselves with an application of a certain subtraction scheme to the QED part of the radiative corrections for the process under consideration. We will use here the $`\overline{\mathrm{MS}}`$scheme Bardeen:1978yd , the DIS scheme can be used as well. This allows to avoid the double counting of the initial quark mass singularities contained in our result for the corrections to the free quark cross section and the ones contained in the corresponding PDF. The latter should be also taken in the same scheme with the same factorization scale.
In fact, using the initial condition for the nonโsinglet NLO QED quark structure function, which coincides with the QCD one with the trivial substitution $`C_F\alpha _sQ_i^2\alpha `$, see Ref. Berends:1987ab , one gets the following expression for the terms to be subtracted from the full calculation with massive quarks:
$`\delta ^{\overline{\mathrm{MS}}}`$ $`=`$ $`{\displaystyle \underset{i=1,2}{}}Q_i^2{\displaystyle \frac{\alpha }{2\pi }}{\displaystyle \underset{0}{\overset{1}{}}}\mathrm{d}\xi _i[{\displaystyle \frac{1+\xi _i^2}{1\xi _i}}(\mathrm{ln}{\displaystyle \frac{M^2}{m_i^2}}`$ (12)
$``$ $`2\mathrm{ln}(1\xi _i)1)]_+\widehat{\sigma }_0(\xi _i),`$
where $`Q_i`$ and $`m_i`$ denote the charge and the mass of the given quark; $`M`$ is the factorization scale; $`\widehat{\sigma }_0(\xi _i)`$ is the cross section at the partonic level with the reduced value of the quark momentum: $`p_i\xi _ip_i`$. The subtracted partonic cross section with $`๐ช\left(\alpha \right)`$ corrections is given by
$`\widehat{\sigma }_1^{\overline{\mathrm{MS}}}=\widehat{\sigma }_1\delta ^{\overline{\mathrm{MS}}}.`$ (13)
Then it can be convoluted with PDFโs as shown below in Eq. (4).
But there is an alternative way to perform the subtraction. Really, to avoid the double counting of the quark mass singularities, we can leave them in the corrected cross section, but remove from the PDFโs:
$`\overline{q}(x,M^2)=q(x,M^2){\displaystyle _x^1}{\displaystyle \frac{\mathrm{d}z}{z}}q({\displaystyle \frac{x}{z}},M^2){\displaystyle \frac{\alpha }{2\pi }}Q_q^2`$
$`\times \left[{\displaystyle \frac{1+z^2}{1z}}\left\{\mathrm{ln}\left({\displaystyle \frac{M^2}{m_q^2}}\right)2\mathrm{ln}(1z)1\right\}\right]_+`$
$`q(x,M^2)\mathrm{\Delta }q,`$ (14)
where $`q(x,M^2)`$ can be taken directly from the existing PDFโs in the $`\overline{\mathrm{MS}}`$scheme (see Ref. Wackeroth:1996hz for the corresponding formula in the DIS scheme). It can be shown analytically (see i.e. Ref. Wackeroth:1996hz ), that this procedure is equivalent to the subtraction from the cross section, and that it really removes (hides) the dependence on the quark masses. The advantage of the last approach is that it can be used regardless of the way to represent the partonic cross section: it can be kept even in the completely differential form.
The natural choices of the factorization scale are $`M^2=M_W^2`$ (when the returning to the $`W`$-resonance is allowed by kinematic cuts) and $`M^2=\widehat{s}=x_1x_2s`$. Variations with respect to the choice should be studied.
In order to avoid the appearance of spurious higher order terms for the case of subtraction from PDFโs, we suggest to apply a procedure of linearization. Schematically it can be represented as follows:
$`\overline{q}_1(x_1,M^2)\times \overline{q}_2(x_2,M^2)\times \widehat{\sigma }_1=[q_1(x_1,M^2)\mathrm{\Delta }q_1]`$
$`\times [q_2(x_2,M^2)\mathrm{\Delta }q_2]\times (\widehat{\sigma }_{\mathrm{Born}}+\widehat{\sigma }_\alpha )`$
$`q_1(x_1,M^2)\times q_2(x_2,M^2)\times \widehat{\sigma }_{\mathrm{Born}}`$
$`+q_1(x_1,M^2)\times q_2(x_2,M^2)\times \widehat{\sigma }_\alpha `$ (15)
$`[q_1(x_1,M^2)\times \mathrm{\Delta }q_2+q_2(x_2,M^2)\times \mathrm{\Delta }q_1]\times \widehat{\sigma }_{\mathrm{Born}},`$
where $`\widehat{\sigma }_{\mathrm{Born}}`$ and $`\widehat{\sigma }_\alpha `$ denote the Born-level partonic cross section and the $`๐ช\left(\alpha \right)`$ RC contribution to it, respectively. Without the linearization procedure, terms with quark mass singularities would remain in the $`๐ช\left(\alpha ^2\right)`$ contribution to the cross section.
## 4 Radiative Corrections to Hadronic processes
The double-differential cross section of the Drell-Yan process can be obtained from the convolution of the partonic cross section with the quark density functions:
$`{\displaystyle \frac{\mathrm{d}\sigma _{\mathrm{RC}}^{pp\mu ^+\nu X}(s)}{\mathrm{d}c\mathrm{d}E_\mu }}={\displaystyle \underset{q_1q_2}{}}{\displaystyle \underset{0}{\overset{1}{}}}{\displaystyle \underset{0}{\overset{1}{}}}dx_1dx_2\overline{q}_1(x_1,M^2)`$
$`\times \overline{q}_2(x_2,M^2){\displaystyle \frac{\mathrm{d}^2\widehat{\sigma }^{q_1q_2\mu ^+\nu }(\widehat{s})}{\mathrm{d}\widehat{c}\mathrm{d}\widehat{E}_\mu }}๐ฅ\mathrm{\Theta }(c,E_\mu ),`$ (16)
where the parton densities with bars mean the ones modified by the subtraction of the quark mass singularities; the step function $`\mathrm{\Theta }(c,E_\mu )`$ defines the phase space domain corresponding to the given event selection procedure. The partonic cross section is taken in the center-of-mass reference frame of the initial quarks, where the cosine of the muon scattering angle, $`\widehat{c}`$, and the muon energy, $`\widehat{E}_\mu `$, are defined. The transformation into the observable variables $`c`$ and $`E_\mu `$ involves the Jacobian:
$`๐ฅ={\displaystyle \frac{\widehat{c}}{c}}{\displaystyle \frac{\widehat{E}_\mu }{E_\mu }}={\displaystyle \frac{4x_1x_2}{a^2}}\sqrt{{\displaystyle \frac{a^2(1+c)}{x_1[a+x_2(1+c)]}}},`$
$`a=x_1+x_2c(x_1x_2),\widehat{c}=1(1c){\displaystyle \frac{2x_1}{a}},`$
$`\widehat{s}=sx_1x_2,\widehat{E}_\mu ={\displaystyle \frac{\sqrt{\widehat{s}}}{2}},`$
$`\widehat{E}_\mu =E_\mu \sqrt{{\displaystyle \frac{1c^2}{1\widehat{c}^2}}}.`$ (17)
An analogous formula can be written for any other choice of a differential distribution as well as for the total cross section.
## 5 Numerical Results and Conclusions
For numerical evaluations we take the same set of input parameters as the one given by Eq. (4.1) of Ref. Dittmaier:2001ay . In Table 2 we present the results for the total cross section<sup>1</sup><sup>1</sup>1factor $`|V_{ud}|^2`$ has been dropped in the sake of comparison with Ref. Dittmaier:2001ay . of the process $`u+\overline{d}\nu _l+l^+(+\gamma )`$. For the Born-level cross section we completely (in all listed digits) agree with the numbers given in Ref. Dittmaier:2001ay . The third line shows radiative corrections in percent before the subtraction of quark mass singularities. These numbers were received directly from the SANC system for the $`G_F`$ EW scheme. Starting from the fourth line we use the treatment of the EW scheme<sup>2</sup><sup>2</sup>2In $`G_F^{}`$ scheme we assigned the following one-loop value of the coupling constant standing at the photon vertices: $`\alpha _{QED}1/132.544`$., which has been adopted Ref. Dittmaier:2001ay . The results for the radiative corrections with $`\overline{\mathrm{MS}}`$subtraction (with factorizations scale being equal to $`M_W`$) are also in a fair agreement. The small deviations there can be due to details in the treatment of EW scheme with respect induced higher order effects. Huge positive corrections in the case without subtraction of quark mass singularities above the $`W`$-peak are due to the initial state radiation which provides the radiative return to the $`W`$-resonance.
The effect of EW scheme dependence is illustrated by Table 3. Results for the total partonic cross section at the Born and $`๐ช\left(\alpha \right)`$ levels are given for two EW schemes. At the Born level the $`7.3\%`$ difference appears just due to the difference in the definition of EW constants in the $`G_F`$ and in the $`\alpha (0)`$ schemes. As it should be the difference between the corrected cross sections is less than the one at the Born level. But still it is large and comparable with the ordered precision of the calculation. Certainly, usage of the $`\alpha (0)`$ is not well motivated for the given energy range. And the difference $`\delta _1`$ gives only an upper estimate of the uncertainty due to the EW scheme dependence. In any case we are going to perform further studies of this effect.
Table 4 represents the dependence of the hadronic DrellโYan cross section on the values of the quark masses with and without the subtraction procedure. The conditions are as follows: the center-of-mass energy is 200 GeV, all events with the invariant mass of the neutrino and charged lepton pair above $`\sqrt{40}`$ GeV are accepted. $`\sigma _0`$ denotes the Born-level cross section obtained using the CTEQ4L set of PDFโs Lai:1996mg . $`\sigma _1`$, $`\sigma _1^{\overline{\mathrm{MS}}(\sigma )}`$, $`\sigma _1^{\overline{\mathrm{MS}}(q)}`$, and $`\sigma _1^{\overline{\mathrm{MS}}(q)}`$(lin.) stand for the cross sections with one-loop EW RC included. The double counting of the quark mass singularities in $`\sigma _1`$ is not removed. The $`\overline{\mathrm{MS}}`$procedure (13) is applied to the partonic cross section in the computation of $`\sigma _1^{\overline{\mathrm{MS}}(\sigma )}`$. Values of $`\sigma _1^{\overline{\mathrm{MS}}(q)}`$ and $`\sigma _1^{\overline{\mathrm{MS}}(q)}`$(lin.) are computed by convolution of the quark (parton) density function modified according to Eq. (3.1) with the full (including quark mass singularities) partonic cross section. The linearization procedure (3.1) was adopted for $`\sigma _1^{\overline{\mathrm{MS}}(q)}`$(lin.) in addition. One can see that the numerical effect of linearization for the given setโup is small (but visible). The two approaches to remove the double counting give very close results as it should be.
For an internal test of our calculations, a comparison of the results produced by our Monte Carlo (MC) and semiโanalytical (SA) codes for the description of hard photon contributions was performed. The results are presented in Table 5, where the corresponding contributions to the protonโproton cross section at 14 TeV center-of-mass energy are given. The conditions and the input parameters were the taken as the ones used in Ref. tuned :
$`\begin{array}{ccccccccc}G_F\hfill & =& 1.16637\times 10^5\mathrm{GeV}^2,\hfill & & & & & & \\ \alpha (0)\hfill & =& 1/137.03599911,\hfill & \alpha _s\hfill & =& 0.1187,\hfill & & & \\ M_W\hfill & =& 80.425\mathrm{GeV},\hfill & \mathrm{\Gamma }_W\hfill & =& 2.124\mathrm{GeV},\hfill & & & \\ M_Z\hfill & =& 91.1867\mathrm{GeV},\hfill & \mathrm{\Gamma }_Z\hfill & =& 2.4952\mathrm{GeV},\hfill & & & \\ M_H\hfill & =& 150\mathrm{GeV},\hfill & m_t\hfill & =& 174.17\mathrm{GeV},\hfill & & & \\ m_u\hfill & =& m_d=66\mathrm{MeV},\hfill & m_c\hfill & =& 1.55\mathrm{GeV},\hfill & & & \\ m_s\hfill & =& 150\mathrm{MeV},\hfill & m_b\hfill & =& 4.5\mathrm{GeV},\hfill & & & \\ |V_{ud}|\hfill & =& |V_{cs}|=0.975,\hfill & |V_{us}|\hfill & =& |V_{cd}|=0.222.\hfill & & & \end{array}`$ (26)
The MRST204QED set Martin:2004dh of PDFโs and the $`G_F`$ EW scheme were used. Six values for the cut on the muon transverse momentum, $`P_T`$, are considered. The cut on the muon rapidity is $`|\eta _l|<1.2`$. The cut on the missing momentum was not imposed, since it canโt be realized in the semi-analytical branch. We also show there the values of the total one-loop EW correction, $`\delta \sigma _{tot}^{MC,SA}`$. The Table 5 shows results for two values of the softโhard photon separator, $`\overline{\omega }`$, and justifies the independence of the total correction on it within the accuracy achieved. The separator is defined in the center-of-mass reference frame of the colliding quarks (partons). We stress, that having a semi-analytical branch of calculations served us as a benchmark and helped a lot to adjust the Monte Carlo code.
In this way with help of the automatized SANC system we calculated the complete one-loop radiative corrections to the charged current Drell-Yan cross section. Our results at the partonic level are in a good agreement with the ones published earlier in Ref. Dittmaier:2001ay . The corresponding computer codes in analytical (FORM) and numerical (FORTRAN) formats are available from SANC SANCwww . They can be used as a part of a more general computer program (like a Monte Carlo event generator) to describe the Drell-Yan process in realistic conditions. Further comparison at the hadronic level with analogous calculations of other groups is in progress tuned .
###### Acknowledgements.
We are grateful to C. Carloni Calame, S. Dittmaier, S. Jadach, M. Krรคmer, G. Montagna, O. Nicrosini W. Placzek, A. Vicini, Z. Was for discussions. Three of us (D.B., L.K. and G.N.) are indebted to the directorate of IFJ, (Cracow, Poland) for a hospitality extended to them in AprilโMay 2005, when a part of this study was completed. This work was supported by the INTAS grant 03-51-4007. One of us (A.A.) thanks for support the grant of the Prezident RF (Scinetific Schols 2027.2003.2) and the RFBR grant 04-02-17192.
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# Untitled Document
hep-th/0506141
Flux-vacua in Two Dimensional String Theory
Juan Maldacena and Nathan Seiberg
School of Natural Sciences
Institute for Advanced Study
Einstein Drive, Princeton, NJ 08540
We analyze the two dimensional type 0 theory with background RR-fluxes. Both the 0A and the 0B theory have two distinct fluxes $`q`$ and $`\stackrel{~}{q}`$. We study these two theories at finite temperature (compactified on a Euclidean circle of radius $`R`$) as a function of the fluxes, the tachyon condensate $`\mu `$ and the radius $`R`$. Surprisingly, the dependence on $`q`$, $`\stackrel{~}{q}`$ and $`\mu `$ is rather simple. The partition function is the absolute value square of a holomorphic function of $`y=|q|+|\stackrel{~}{q}|+i\sqrt{2\alpha ^{}}\mu `$ (up to a simple but interesting correction). As expected, the 0A and the 0B answers are related by T-duality. Our answers are derived using the exact matrix models description of these systems and are interpreted in the low energy spacetime Lagrangian.
6/05
1. Introduction
The renewed interest in noncritical string theories has originated from their relevance to current topics in string theory \[1--4\], like open/closed duality, holography and D-branes. These models are interesting because they have a complete nonperturbative definition and at the same time can be analyzed exactly. As such, they are good laboratories for subtle nonperturbative questions. In particular, this is the only framework where flux vacua โ backgrounds with RR-fluxes โ can be analyzed exactly. Issues associated with such flux vacua have already been discussed both in $`\widehat{c}=1`$ models \[3--18\] and in $`\widehat{c}<1`$ models \[19--21\]. However, some of the results of the $`\widehat{c}=1`$ system appeared confusing and it has been suggested that the system with RR-flux is related to black holes. The purpose of this paper is to clarify some of these confusions.
The $`\widehat{c}=1`$ model is a two dimensional string theory. The target space is parameterized by the time $`t`$ and the spatial coordinate $`\varphi `$. The background is not translational invariant; the system has a linear dilaton which makes the string coupling space dependent
$$g_s(\varphi )=e^\varphi $$
The $`\varphi +\mathrm{}`$ asymptotic region is characterized by weak coupling. Scattering experiments are performed by sending signals from this asymptotic region and detecting them as they return. More precisely, the scattering is to and from null infinities $`๐ฅ^\pm `$; the incoming modes are functions of $`\varphi +t`$ and the outgoing modes are functions of $`\varphi t`$. It has been assumed that the system does not have another asymptotic region with $`\varphi \mathrm{}`$; i.e. there is no separate scattering to and from that region. Indeed, the strong coupling region has effectively finite volume (see, e.g. ).
In section 2 we discuss the spacetime picture of the two kinds of models we study, the 0B and the 0A theories. We review their spectra and the leading order terms in the spacetime effective Lagrangian. We show that the 0B theory has two kinds of continuous RR-fluxes, $`\nu `$ and $`\stackrel{~}{\nu }`$, and 0A theory has two kinds of quantized RR-fluxes, $`q`$ and $`\stackrel{~}{q}`$.
In addition to these two parameters we can also turn on a โtachyonโ condensate $`T(\varphi )=\mu e^\varphi `$ and study the physics as a function of the real parameter $`\mu `$. The analysis of showed that the physics is smooth as a function of $`\mu `$. Finally, we can also study the system with Euclidean time which is compactified on a circle of radius $`R`$. This corresponds to studying the thermodynamics of the theory with finite temperature $`\frac{1}{2\pi R}`$.
In section 3 we study the exact 0A theory with its two RR-fluxes and derive an expression for its partition function $`๐ต_A(\mu ,q,\stackrel{~}{q},R)`$ as a function of all variables. We find that up to a simple (but interesting) term, the dependence on $`q`$ and $`\stackrel{~}{q}`$ is only through $`\widehat{q}=|q|+|\stackrel{~}{q}|`$. Furthermore, up to the same simple term, the partition function factorizes as a holomorphic function of $`y=\widehat{q}+i\sqrt{2\alpha ^{}}\mu `$ and its complex conjugate. We interpret the dependence on $`\widehat{q}=|q|+|\stackrel{~}{q}|`$ as a result of the presence of $`|q\stackrel{~}{q}|`$ fundamental strings in the system. This is reminiscent of the factorization involved in topological string computations, see e.g. and references therein.
The presence of two distinct fluxes in the 0A theory, which couple differently to the tachyon, has led to speculations about the existence of extremal black hole solutions when the tachyon is not excited: $`\mu =0`$. Indeed the lowest order in $`\alpha ^{}`$ equations of motion predict such a solution . Unfortunately it is not possible to trust the leading order equations in the two dimensional string theory. The fact that our matrix model results depend only on $`\widehat{q}=|q|+|\stackrel{~}{q}|`$ shows that the physics is essentially the same as the physics with only one kind of flux that has been analyzed in , see also . Such analysis does not show any indications of a black hole, namely there is no entropy and there is no classical absorption. So the matrix model, as analyzed in this paper is not consistent with an object that could be called a black hole.
In section 4 we study the exact 0B theory. We view the partition function of the noncompact Lorentzian theory $`๐ต_B`$ as a transition amplitude between the past and the future. The value of this amplitude is complex, but is simpler than expected. Its phase is given by the real part of $`\mathrm{\Xi }\left(2i\sqrt{2\alpha ^{}}(|\nu |+|\stackrel{~}{\nu }|+\mu )\right)+\mathrm{\Xi }\left(2i\sqrt{2\alpha ^{}}(|\nu |+|\stackrel{~}{\nu }|\mu )\right)`$ for some function $`\mathrm{\Xi }`$ and $`\nu ,\stackrel{~}{\nu }`$ are the Lorentzian RR fluxes. The finite temperature version of the 0B theory has quantized fluxes $`|q|=2iR|\nu |`$ and $`|\stackrel{~}{q}|=2iR|\stackrel{~}{\nu }|`$. The expression for the finite temperature partition function is related to that of the 0A theory by the expected T-duality with the following change in the parameters
$$R_B=\frac{\alpha ^{}}{R_A},\mu _B=\frac{R_A}{\sqrt{2\alpha ^{}}}\mu _A$$
with the same $`q`$ and $`\stackrel{~}{q}`$.
In Appendix A we review and extend a simple formalism for describing these systems . It allows us to simply compute the transition amplitudes both of the 0B and the 0A theory and to obtain new insights into the nature of the scattering.
2. Spacetime effective Lagrangian
In this section we focus on the weak coupling end of the target space, $`\varphi +\mathrm{}`$ and study the low energy field theory there. Since the string coupling is arbitrarily small, the dynamics is dominated by classical physics, and the leading approximation to the effective field theory is valid. In particular, the massless modes are reliably found by a weak coupling worldsheet analysis. Another simplification in this part of the target space is that a possible tachyon condensate
$$T(\varphi )=\mu e^\varphi $$
can be neglected there.
2.1. 0B
We start by analyzing the 0B string theory. The spectrum consists of two massless scalars an NS-NS โtachyonโ $`T`$ and an RR scalar $`C`$.
It is clear from the worldsheet description that the theory has two discrete $`๐_2`$ symmetries :
1. The first symmetry acts in the worldsheet description as $`(1)^{F_L}`$ where $`F_L`$ is the target space fermion number of the worldsheet left movers. It acts on the spectrum as
$$\begin{array}{cc}& TT\hfill \\ & CC\hfill \end{array}$$
Since it changes the sign of the RR scalar, it changes the charge of D-branes; we will can refer to it as charge conjugation.
2. A more subtle symmetry acts in the worldsheet description as $`(1)^{f_L}`$ where $`f_L`$ is the left moving worldsheet fermion number. It acts on the spectrum as
$$\begin{array}{cc}& TT\hfill \\ & C_LC_L\hfill \\ & C_RC_R\hfill \end{array}$$
Here $`C_L`$ and $`C_R`$ are the target space left and right moving components of $`C`$. Hence the action of this symmetry on $`C`$ is a duality transformation. By analogy with its higher dimensional counterpart, we will refer to it as S-duality.
The invariance under S-duality means that the scalar $`C`$ is compact and its radius is the selfdual radius. Therefore, the asymptotic theory as $`\varphi +\mathrm{}`$ has an $`SU(2)\times SU(2)`$ symmetry. This symmetry will be important below.
Since $`T`$ is odd under the S-duality symmetry, the kinetic term of $`C`$ has to be of the form $`\frac{1}{8\pi }f(T)(C)^2`$ with $`f(T)=\frac{1}{f(T)}`$. More detailed worldsheet considerations show that $`f(T)=e^{2T}`$ and hence the kinetic term is
$$_{kinetic}=\frac{1}{8\pi }e^{2T}\left[(_tC)^2(_\varphi C)^2\right]$$
Therefore, the coupling of $`T`$ to $`C`$ breaks the $`SU(2)\times SU(2)`$ symmetry to $`U(1)\times U(1)`$. In particular, the tachyon condensate $`T(\varphi )=\mu e^\varphi `$ breaks the S-duality symmetry; more precisely, the theory with $`\mu `$ is related by S-duality to the theory with $`\mu `$.
Let us examine the equations of motion which arise from (2.1)
$$_t(e^{2T}_tC)_\varphi (e^{2T}_\varphi C)=0$$
For $`\varphi +\mathrm{}`$ we can neglect $`T`$ in this expression and $`C`$ is simply a free scalar at the selfdual radius. Of particular interest to us will be the zero momentum solutions of the equations of motion
$$\frac{C}{\sqrt{2}}2(\nu \varphi +\stackrel{~}{\nu }t)=\nu _{in}(\varphi +t)+\nu _{out}(\varphi t)$$
where we used an approximate sign to remind us that this solution is valid only for $`\varphi +\mathrm{}`$. The two integration constants $`\nu `$ and $`\stackrel{~}{\nu }`$, or equivalently $`\nu _{in}`$ and $`\nu _{out}`$ are RR-fluxes. Both of them are odd under the charge conjugation symmetry (2.1) and transform under S-duality (2.1) as $`\nu \stackrel{~}{\nu }`$, or equivalently, $`\nu _{out}\nu _{out}`$. Since both the $`\nu `$ and $`\stackrel{~}{\nu }`$ deformations are non-normalizable as $`\varphi +\mathrm{}`$ they label backgrounds which are determined by the behavior at infinity and they do not fluctuate.
The coupling to $`T`$ in (2.1) has important consequences. If $`T(\varphi )`$ is nonzero the solution (2.1) becomes
$$\frac{C}{2\sqrt{2}}=\stackrel{~}{\nu }t+\nu e^{2T(\varphi )}๐\varphi $$
(note, as a check that as $`\varphi +\mathrm{}`$ it goes over to (2.1)). Consider the effect of $`T(\varphi )=\mu e^\varphi `$ with positive $`\mu `$ on (2.1). At the strong coupling end $`\varphi \mathrm{}`$ the second term $`\nu e^{2T(\varphi )}๐\varphi `$ rapidly goes to zero, and the corresponding mode is normalizable. (Of course, it is not normalizable as $`\varphi +\mathrm{}`$.) Hence it is a standard background deformation.
This is to be contrasted with the first term $`\stackrel{~}{\nu }t`$. The norm of the small fluctuations which is derived from (2.1) is $`๐\varphi e^{2T(\varphi )}\delta C^2`$, and hence it is not normalizable at $`\varphi \mathrm{}`$. Such a deformation, which is singular in the interior of the target space, can be present only if an object is present at its singularity. In our case, the relevant object is a D-brane which carries RR-charge. It sources the RR-flux $`\stackrel{~}{\nu }`$.
For negative $`\mu `$ the situation is reversed. The RR-flux $`\stackrel{~}{\nu }`$ is normalizable at $`\varphi \mathrm{}`$, and it does not need a D-brane source. However, the other flux $`\nu e^{2T(\varphi )}๐\varphi `$ needs D-branes at $`\mathrm{}`$. This exchange in the behavior of the two fluxes under the change of the sign of $`\mu `$ is consistent with the S-duality symmetry.
As we vary $`\mu `$ from positive to negative values the physics has to change in a continuous fashion. This is particularly obvious in the asymptotic weak coupling end where the effects of nonzero $`\mu `$ are negligible. Therefore, we see here a phenomenon which has already been observed elsewhere \[26--30,,19--21\], that RR-flux without D-branes can be continuously transformed to RR-flux carried by D-branes.
We should clarify the nature of these D-branes at infinity. In the worldsheet description these are the so called ZZ-branes . Since they couple to the scalar $`C`$, the relevant branes are D-instantons. This means that our background flux represents a transition which is mediated by instantons. For positive $`\mu `$ we have $`\stackrel{~}{\nu }`$ D-instantons per unit time and for negative $`\mu `$ we need $`\nu `$ such instantons per unit time. Although the number of such D-instantons is quantized, the numbers per unit time, $`\nu `$ or $`\stackrel{~}{\nu }`$ do not have to be quantized.
Let us examine a background with generic $`\nu `$ and $`\stackrel{~}{\nu }`$. It is easy to calculate the energy momentum tensor of that background as $`\varphi +\mathrm{}`$. It is
$$T_{++}=\frac{1}{4\pi }\nu _{in}^2+\mathrm{},T_{}=\frac{1}{4\pi }\nu _{out}^2+\mathrm{},T_+=0+\mathrm{}$$
Here the ellipses represent $`\varphi `$ dependent corrections which are negligible as $`\varphi +\mathrm{}`$. As far as the asymptotic Lagrangian (2.1) is concerned, there is no problem with such a background. However, a crucial part of the story is that the dynamics is such that pulses get reflected from the $`\varphi =\mathrm{}`$ region. Furthermore the reflection from this region conserves energy. On the other hand, the incoming energy flux from $`๐ฅ^{}`$ which is $`๐x^{}T_{}`$ is not the same as the outgoing energy flux through $`๐ฅ^+`$ which is $`๐x^+T_{++}`$. (These integrals are infinite since we have a constant flux). Therefore, conservation of energy implies that we should either send in extra excitations from the past, or produce extra excitations in the future. For simplicity we can focus on the lowest energy excitation by adding excitations on the side that has the lower flux, so as to match the side with higher flux. The lowest energy configuration has
$$T_{++}=T_{}=\frac{1}{4\pi }\mathrm{max}(\nu _{in}^2,\nu _{out}^2)+\mathrm{}=\frac{1}{4\pi }(|\nu |+|\stackrel{~}{\nu }|)^2+\mathrm{}$$
Depending on whether $`\nu _{out}^2`$ is bigger or smaller than $`\nu _{out}^2`$, this is achieved by adding waves with $`T_{t\varphi }=\frac{1}{\pi }\nu \stackrel{~}{\nu }+\mathrm{}`$ either in the past or in the future .
We are going to be interested in computing the scattering amplitude between a state in the past which is characterized by $`\nu _{in}`$ and a state in the future which is characterized by $`\nu _{out}`$. A full characterization of the states involves specifying the state for all the oscillators of the fields $`T`$ and $`C`$. All we are doing in this section is to analyze the asymptotic region in order to understand which states we can send in and which states we expect to come out. Below, we will extend this discussion in the asymptotic region to the full system and will derive the exact expression for for the scattering amplitude. We will see that it depends only on $`|\nu |+|\stackrel{~}{\nu }|`$.
Finally, we would like to comment on the 0B theory on a Euclidean circle of radius $`R`$. The analytic continuation to Euclidean space leads to real $`\nu _E=i\nu `$ and $`\stackrel{~}{\nu }_E=i\stackrel{~}{\nu }`$, where the subscript $`E`$ denotes that these are the Euclidean values. Here, in this Euclidean time setup for positive $`\mu `$ the parameter $`\nu _E`$ is proportional to the number of instantons per unit Euclidean time and therefore the parameter $`q=2\nu _ER`$ is quantized (the precise normalization will be derived below). Similarly, for negative $`\mu `$ the parameter $`\stackrel{~}{q}=2\stackrel{~}{\nu }_ER`$ is quantized. By continuity these two parameters are quantized for all $`\mu `$.
2.2. 0A
The discussion of the 0A string theory parallels that of the 0B theory. In fact, when the 0B theory is analytically continued to Euclidean time and that coordinate is compactified, it is T-dual to the 0A theory.
The spectrum of the 0A theory consists of an NS-NS โtachyonโ $`T`$, but the RR-scalar $`C`$ is absent. It is replaced by two gauge fields $`F_{t\varphi }`$ and $`\stackrel{~}{F}_{t\varphi }`$. These gauge fields have no propagating degrees of freedom, and only their zero momentum values can change.
Again, the theory has two discrete $`๐_2`$ symmetries:
1. The charge conjugation symmetry which acts on the worldsheet theory as $`(1)^{F_L}`$ acts on these fields as
$$\begin{array}{cc}& TT\hfill \\ & F_{t\varphi }F_{t\varphi }\hfill \\ & \stackrel{~}{F}_{t\varphi }\stackrel{~}{F}_{t\varphi }\hfill \end{array}$$
2. The S-duality symmetry which acts on the worldsheet theory as $`(1)^{f_L}`$ acts on the fields as
$$\begin{array}{cc}& TT\hfill \\ & F_{t\varphi }\stackrel{~}{F}_{t\varphi }\hfill \\ & \stackrel{~}{F}_{t\varphi }F_{t\varphi }\hfill \end{array}$$
The Lagrangian (2.1) is replaced by
$$=\pi \alpha ^{}\left(e^{2T}F_{t\varphi }^2+e^{2T}\stackrel{~}{F}_{t\varphi }^2\right)$$
which is invariant under the two symmetries (2.1)(2.1). The asymptotic solution of the equations of motion is $`2\pi \alpha ^{}F_{t\varphi }=q`$, $`2\pi \alpha ^{}\stackrel{~}{F}_{t\varphi }=\stackrel{~}{q}`$. Including the $`T`$ dependent prefactors in (2.1) the solutions are
$$\begin{array}{cc}& F_{t\varphi }=qe^{2T}\hfill \\ & \stackrel{~}{F}_{t\varphi }=\stackrel{~}{q}e^{2T}\hfill \end{array}$$
For negative $`\mu `$ background $`F_{t\varphi }`$ is singular at $`\varphi \mathrm{}`$ and is generated by D-branes at $`\varphi \mathrm{}`$, while the background $`\stackrel{~}{F}_{t\varphi }`$ is regular and does not need such branes. For positive $`\mu `$ the situation is reversed. These D-branes at infinity carry RR-electric charge. As in the 0B theory, these are charged ZZ-branes. However, unlike the D-instantons of the 0B theory the relevant branes couple to gauge field one forms, and therefore they are D0-branes. Hence, $`q`$ and $`\stackrel{~}{q}`$ are quantized.
Consider a background with generic quantized values of $`q`$ and $`\stackrel{~}{q}`$. By analogy with similar situations in critical string theory \[32--35\] we expect that such a background is possible only if we add to the system $`q\stackrel{~}{q}`$ fundamental strings. This expectation can be derived by examining the coupling to the two form field $`B`$. Such a field does not have interesting dynamics in two dimensions, but its equation of motion shows that such strings must be present. This conclusion can also be derived by starting with the Euclidean 0B theory on a circle. In the 0B theory we had to add energy flux to compensate the imbalance $`T_{t\varphi }\frac{1}{\pi }\nu \stackrel{~}{\nu }`$. This translates, after rotation to Euclidean space and T-duality, to adding $`q\stackrel{~}{q}`$ strings in the 0A theory. Note that the sign of $`q\stackrel{~}{q}`$ is correlated with the orientation of these strings in the two dimensional target space.
3. 0A Matrix model
In this section we consider the matrix model of the two dimensional 0A string theory . This is a gauged matrix model which contains a complex matrix, $`m`$, which transforms in the bifundamental of $`U(N)_A\times U(N)_B`$. There are two ways of introducing fluxes. First, we can modify the gauge groups so that we start with $`U(N)_A\times U(M)_B`$ with $`q=MN`$ (for simplicity assume that $`q>0`$). This leads to $`\stackrel{~}{q}=0`$, $`q0`$. This corresponds to placing $`M`$ charged ZZ branes and $`N`$ anti-ZZ-branes at $`\varphi =\mathrm{}`$ and then letting the open string tachyon condense. It is clear from this description that as long as $`\mu `$ is below the barrier (in our conventions, $`\mu <0`$), we will have $`q`$ charged ZZ branes left over. So in this set up we end up describing the configuration with the flux that is sourced by D-branes. We expect that these ZZ branes will be stuck at the strong coupling end since a charged ZZ brane does not have an open string tachyon. One could consider charged D0 branes that move in the bulk of the two dimensional space . We expect that these D0 branes will exist only in the non-singlet sector of the matrix model, since the Euclidean boundary states that represent them contain a non-normalizable open string winding mode.
The second way to introduce flux is for $`\stackrel{~}{q}0`$, $`q=0`$. In this case we set $`N=M`$ and add to the matrix model a term of the form
$$S=S_0+i\stackrel{~}{q}(TrATrB)$$
where $`A`$ and $`B`$ are the gauge fields for the $`U(N)_A`$ and $`U(N)_B`$ gauge groups respectively. As shown in this leads to a problem where the eigenvalues move in a complex plane, all with angular momentum $`\stackrel{~}{q}`$. This can also be viewed as a special case of the general problem of coupling the matrix model to non-singlet representations. In this case we simply have a singlet representation of $`SU(N)_A\times SU(N)_B`$ which carries charge under the relative $`U(1)`$ (which is generated by the difference between the generators of $`U(1)_AU(N)_A`$ and of $`U(1)_BU(N)_B`$). Below the barrier, $`\mu <0`$, we can understand the origin of (3.1) as follows. As we explained above, the charged ZZ branes source the flux $`F`$. In this case the second flux $`\stackrel{~}{F}`$ can be excited and leads to a smooth geometry. If we add ZZ branes the flux $`\stackrel{~}{F}`$ leads to a Chern-Simons term on the worldvolume of the ZZ branes that produces (3.1). This is the same type of coupling that leads to the usual Chern-Simons terms on D-brane worldvolumes in the ten dimensional superstring.
Surprisingly, once we reduce the problem to eigenvalues, the dynamics of these two cases is exactly the same . Below we will slightly qualify this general comment.
We can now study the case with non-zero $`q`$ and $`\stackrel{~}{q}`$. The first naive idea is to consider again a rectangular matrix with $`M=N+q`$ and add the Chern-Simons term (3.1). Let us assume for simplicity that $`q>0`$. However, as we now explain, in this case the path integral vanishes. The matrix model degree of freedom, the matrix $`m`$, is not charged under the the diagonal $`U(1)`$ generated by the sum of the generators of $`U(1)_A`$ and $`U(1)_B`$. Our normalization for these $`U(1)`$s is such that the fundamental representation of $`SU(N)U(N)`$ has charge one (modulo $`N`$). On the other hand, the coupling (3.1) leads to charge $`q\stackrel{~}{q}`$ under diagonal $`U(1)`$. Since this charge cannot be cancelled, the path integral vanishes. In other words we cannot obey the Gauss law for this gauge field.
In order to learn how to deal with this, let us return to the spacetime picture and understand what happens in string theory when we start with flux $`\stackrel{~}{q}`$ and we attempt to put a charged ZZ brane. There is a coupling $`S=i\stackrel{~}{q}B`$ on the worldvolume of the ZZ brane, where $`B`$ is the worldvolume $`U(1)`$ gauge field. In order to cancel this charge we need to have $`\stackrel{~}{q}`$ strings ending on the ZZ brane. If we have $`q`$ charged ZZ branes we need to add $`q\stackrel{~}{q}`$ strings ending on them. A similar situation has been encountered in various situations in \[32--35\]. So the matrix model that contains both fluxes necessarily involves the presence of a non-trivial representation of $`U(M)_B`$ with $`M`$-ality $`q\stackrel{~}{q}`$.
In summary, the matrix model with both fluxes is a $`U(N)_A\times U(N+q)_B`$ gauged matrix model
$$๐(A,B,m)e^{i{\scriptscriptstyle ๐tTr[(D_0m)^{}D_0m+{\scriptscriptstyle \frac{1}{2\alpha ^{}}}m^{}m]}}e^{i\stackrel{~}{q}{\scriptscriptstyle (TrATrB)}}Tr_{}Pe^{i{\scriptscriptstyle B}}$$
where $``$ is a representation of $`U(N+q)_B`$ with $`q\stackrel{~}{q}`$ $`M`$-ality<sup>1</sup> In principle we can also introduce a representation of $`U(N)_A`$. In this case the constraint is that the $`M`$-ality of the representation of $`U(M)_B`$ minus the $`N`$-ality of the representation of $`U(N)_A`$ should be $`q\stackrel{~}{q}`$.. We can now analyze this problem using the general procedure described in . We diagonalize the matrix $`m`$ and we integrate out the gauge fields. Then we get an effective hamiltonian of the form<sup>2</sup> The apparent differences between this expression and the one in are due to the fact that in $`\stackrel{~}{q}`$ is included as part of the $`U(1)`$ charge of the representation.
$$\begin{array}{cc}\hfill H=& [\underset{i=1}{\overset{N}{}}\frac{1}{2}\frac{^2}{\rho _i}\frac{1}{2}\rho _i^2+\frac{1}{2}\frac{\stackrel{~}{q}^2+q^2\frac{1}{4}}{\rho _i^2}+\hfill \\ \hfill +& 2\underset{i<jN}{}\frac{\mathrm{\Pi }_j^i\mathrm{\Pi }_i^j}{(\rho _i^2\rho _j^2)}+\underset{i=1}{\overset{N}{}}\frac{1}{\rho _i^2}\underset{j>N}{}(\mathrm{\Pi }_j^i\mathrm{\Pi }_i^j+\mathrm{\Pi }_i^j\mathrm{\Pi }_j^i)]P_0\hfill \end{array}$$
where $`\mathrm{\Pi }_i^j`$ are the $`U(N+q)`$ generators in the representation $``$ and $`P_0`$ is a projector on the states obeying
$$\begin{array}{cc}& \mathrm{\Pi }_i^i=0(\mathrm{no}\mathrm{sum}),i=1,\mathrm{},N\hfill \\ & \mathrm{\Pi }_l^k=\stackrel{~}{q}\delta _l^k,l,k>N\hfill \end{array}$$
The last condition implies that under the decomposition $`U(N+q)U(N)\times U(q)`$ we select states that are singlets of $`SU(q)`$ and carry $`U(1)_qU(q)`$ charge $`q\stackrel{~}{q}`$. The simplest way in which we can achieve this is by starting out with an $`SU(N+q)`$ representation whose Young tableaux contains $`q`$ rows of length $`\stackrel{~}{q}`$ (we assume that $`\stackrel{~}{q}>0`$), see figure 1(a). Let us call this representation $`_0`$. In this case the state that is a singlet under $`SU(q)`$ is also a singlet under $`SU(N)`$ and obeys the two conditions (3.1).
Fig. 1: The Young tableaux (a) corresponds to the simplest representation $`_0`$ which leads to a nontrivial answer. The Young tableaux (b) is a more complicated representation which also contributes. The letters and numbers along the sides of the diagrams denote the number of boxes in that side.
We now need to consider the operator that appears in the Hamiltonian (3.1)
$$Q_i=\underset{j>N}{}(\mathrm{\Pi }_j^i\mathrm{\Pi }_i^j+\mathrm{\Pi }_i^j\mathrm{\Pi }_j^i),iN$$
This operator transforms in the singlet of $`SU(q)`$ and it decomposes as the singlet plus adjoint in $`SU(N)`$. This implies that when we act on the single state that is $`SU(q)`$ invariant in the representation $`_0`$ it can give us a state in the adjoint or the singlet of $`SU(N)`$. Since the only state in $`_0`$ that is in the singlet of $`SU(q)`$ is also in the singlet of $`SU(N)`$, we conclude that the action of $`Q_i`$ can only give us a singlet. So this action will just be a c-number. We can simply compute this c-number to be
$$Q_i=q\stackrel{~}{q}$$
Going back to the hamiltonian we find that it reduces to
$$H=\underset{i=1}{\overset{N}{}}\frac{1}{2}\frac{^2}{\rho _i}\frac{1}{2}\rho _i^2+\frac{1}{2}\frac{(\stackrel{~}{q}+q)^2\frac{1}{4}}{\rho _i^2}$$
Here we have assumed that $`q`$ and $`\stackrel{~}{q}`$ are positive. The same analysis can be repeated for the general case and we find that the dynamics depends only on
$$\widehat{q}=|q|+|\stackrel{~}{q}|$$
This is a surprising result from the point of view of the target space theory, as well as the matrix model.
Before we continue, let us explain what happens if our representation is a more general representation that contains a state obeying (3.1). An example of a more general representation can be found in figure 1(b). In this case a state that obeys the second condition in (3.1) can be in the singlet or the adjoint of $`SU(N)`$. The operator in (3.1) mixes the singlet with the adjoint of $`SU(N)`$. So we expect that all these states will have properties that are similar to those encountered in general non-singlet representations as discussed in . Such states have a divergent energy gap compared to the state that comes from the representation $`_0`$ in figure 1(a). In the spacetime theory these states can be understood as states that, besides the $`q\stackrel{~}{q}`$ strings ending on the charged ZZ brane, contain more string anti-string pairs. The fact that these strings stretch all the way to infinity is related to this divergence in the energy . This divergence implies that these other states are in a different superselection sector. We would have found similar divergencies, due to extra strings, if we had also introduced a representation of the first group $`U(N)_A`$. So from now on we will assume that we are adding simply the representation $`_0`$ in figure 1(a).
Finally, let us present an alternate physical interpretation of the need for the representation $`_0`$. When the two kinds of branes/fluxes are present the system has massive open string fermions<sup>3</sup> These are similar to the fermions in the D0-D8 system, except that here they come from open strings stretched all the way to infinity, and hence they are infinitely massive.. These can be added to the matrix model. Their quantization leads to several representations which ultimately, after the use of the Gauss law constraints, lead to $`_0`$.
The finite temperature partition function can then be computed as in where the case $`\stackrel{~}{q}=0`$ and arbitrary $`q`$ was studied. We repeat this computation in Appendix A. Our arguments that it is a function of $`\widehat{q}=|q|+|\stackrel{~}{q}|`$ alow us to extend it to
$$_\mu ^3\mathrm{log}๐ต_{0A}=Re\frac{1}{2}_0^{\mathrm{}}\frac{dt}{t}_\mu ^3e^{(\widehat{q}+i\sqrt{2\alpha ^{}}\mu )\frac{t}{2}}\frac{1}{\mathrm{sinh}\frac{t}{2}\mathrm{sinh}\sqrt{\frac{\alpha ^{}}{2}}\frac{t}{2R}}$$
Note that the integral in the right hand side converges. When we integrate this expression three times with respect to $`\mu `$ we need three integration constants โ a $`q`$ dependent second order polynomial in $`\mu `$. We claim that the answer is
$$\begin{array}{cc}& \mathrm{log}๐ต_{0A}(\mu ,q,\stackrel{~}{q},R)=\mathrm{\Omega }(y,r)+\mathrm{\Omega }(\overline{y},r)+(2\pi R)\frac{\mu }{4}(|q||\stackrel{~}{q}|)\hfill \\ & y=|q|+|\stackrel{~}{q}|+i\sqrt{2\alpha ^{}}\mu \hfill \\ & r=R\sqrt{\frac{2}{\alpha ^{}}}\hfill \end{array}$$
where the function $`\mathrm{\Omega }(y,r)`$ is given by
$$\mathrm{\Omega }(y,r)_0^{\mathrm{}}\frac{dt}{t}\left[e^{\frac{yt}{2}}\frac{1}{4\mathrm{sinh}\frac{t}{2}\mathrm{sinh}\frac{t}{2r}}\frac{r}{t^2}+\frac{ry}{2t}+[\frac{1}{24}(r+\frac{1}{r})\frac{ry^2}{8}]e^t\right]$$
We are interested in $`Re(y)=\widehat{q}0`$ where this integral converges. (More generally, it converges for $`Re(y)>(1+\frac{1}{r})`$.) Therefore, $`\mathrm{\Omega }`$ is an analytic function of $`y`$ which is real along the positive real $`y`$ axis. It is interesting that a closely related function appears in a totally different context in <sup>4</sup> In terms of the function $`G(x)`$ in our function is $`\mathrm{\Omega }(y,r)=C(b)+\mathrm{log}G(\frac{y}{2b}+\frac{Q}{2})`$ where $`b=\sqrt{r}`$, $`Q=b+b^1`$ and $`C(b)`$ is a constant independent of $`y`$..
The expression (3.1) for the partition function is one of the main results of this paper. Let us discuss it in more detail.
The last term in (3.1) depends on $`q`$ and $`\stackrel{~}{q}`$ separately and not only on $`\widehat{q}=|q|+|\stackrel{~}{q}|`$. We will return to it below.
It is surprising that up to this last term in (3.1) the complicated function $`\mathrm{log}๐ต_{0A}`$ is given as a sum of a holomorphic and an anti-holomorphic functions. Correspondingly, the partition function $`๐ต_{0A}`$ satisfies holomorphic factorization. The polynomial in $`\mu `$ which is not determined by (3.1) was fixed such that this holomorphic factorization is satisfied. In the next section we will present another computation of this partition function where some of this polynomial dependence on $`\mu `$ is independently determined.
It is straightforward to work out the asymptotic expansion of $`\mathrm{\Omega }(y,r)`$ at large $`y`$ with $`Re(y)0`$
$$\mathrm{\Omega }(y,r)=(\mathrm{log}\frac{y}{2}\frac{3}{2})\frac{ry^2}{8}\frac{1}{24}(r+\frac{1}{r})\mathrm{log}\frac{y}{2}\frac{7r^2+10+\frac{7}{r^2}}{1440}\frac{1}{ry^2}+๐ช(\frac{1}{y^4})$$
which leads to the following expression at large $`\mu `$
$$\begin{array}{cc}\hfill \mathrm{\Omega }(y=\widehat{q}+i\sqrt{2\alpha ^{}}\mu ,r)+& \mathrm{\Omega }(\overline{y}=\widehat{q}i\sqrt{2\alpha ^{}}\mu ,r)=\left(\frac{3}{2}\mathrm{log}(\sqrt{\frac{\alpha ^{}}{2}}|\mu |)\right)\frac{\alpha ^{}r\mu ^2}{2}\frac{2\pi R|\mu |\widehat{q}}{4}\hfill \\ & +\left(\frac{\widehat{q}^2r}{4}\frac{1}{12}(r+\frac{1}{r})\right)\mathrm{log}(\sqrt{\frac{\alpha ^{}}{2}}|\mu |)\hfill \\ & +\frac{\frac{7}{r^2}+10+7r^2+15\widehat{q}^2r(\widehat{q}^2r2(r+\frac{1}{r}))}{1440\alpha ^{}r\mu ^2}+๐ช\left(\frac{1}{\mu ^4}\right)\hfill \end{array}$$
In the worldsheet genus expansion these terms have the following interpretation. The first term corresponds to the sphere contribution. The scaling of the second term which is proportional to $`\widehat{q}`$ shows that it arises from a disk diagram. We will return to this term below. The $`\widehat{q}`$ independent term in the second line is the contribution of the torus and the term proportional to $`\widehat{q}^2`$ corresponds to a sphere with two RR insertions, or an annulus. Higher orders can be discussed similarly.
Let us focus on the term $`2\pi R|\mu |\widehat{q}/4`$ in (3.1), which comes from the first term in (3.1). Despite appearance, because of the absolute value sign on $`\mu `$, it should not be discarded as an un-interesting analytic term. We can now understand the role of the last term in (3.1). Combining it with $`2\pi R|\mu |\widehat{q}/4`$ we have
$$\frac{2\pi R|\mu |}{4}(|q|+|\stackrel{~}{q}|)+\frac{2\pi R\mu }{4}(|q||\stackrel{~}{q}|)=\{\begin{array}{cc}\pi R\mu |\stackrel{~}{q}|\hfill & \mu >0\hfill \\ \pi R|\mu ||q|\hfill & \mu <0\hfill \end{array}$$
We interpret this as the contribution of the disk amplitude of the charged ZZ-branes. For positive $`\mu `$ we have $`|\stackrel{~}{q}|`$ ZZ-branes and for negative $`\mu `$ we have $`|q|`$ ZZ-branes. Each has energy $`|\mu |/2`$. (Recall that the energy of a brane-anti-brane pair is equal to $`|\mu |`$. So the energy of a single charged D-brane should be equal to $`|\mu |/2`$.). From the point of view of the 0A matrix model we needed to introduce the last term in (3.1) โby handโ in order to obey (3.1). This is an analytic term in $`\mu `$, so it is, in principle, possible to introduce it. However, we will see that this term emerges naturally from the 0B matrix model.
So after including the analytic term we find precisely the expected behavior for the free energy. For $`\mu <0`$ we have D-branes that produce flux proportional to $`q`$ and there are no terms that have odd powers in $`\stackrel{~}{q}`$ in the asymptotic expansion. On the other hand for $`\mu >0`$ we have the opposite situation, since now the flux $`\stackrel{~}{q}`$ is sourced by D-branes.
Despite this simple physical interpretation, our result is still surprising. With the exception of the disk term (3.1) the semiclassical expansion includes only even powers of $`\mu `$ and $`q`$ and $`\stackrel{~}{q}`$. This means that there are no contributions from worldsheets with odd number of boundaries. For positive $`\mu `$ this is the expected result when $`\stackrel{~}{q}=0`$ and there are no ZZ-branes. Similarly, for negative $`\mu `$ this is the expected result when $`q=0`$. The dependence on $`q`$ and $`\stackrel{~}{q}`$ through $`\widehat{q}=|q|+|\stackrel{~}{q}|`$ together with these expected results guarantee that, with the exception of the disk (3.1), there are no contributions from surfaces with odd number of boundaries. We do not have a worldsheet or spacetime interpretation of this surprising result.
4. 0B Matrix model
4.1. Lorentzian 0B model
In this section we consider the 0B matrix model which consists of a hermitian matrix model with an inverted harmonic oscillator potential such that in the free fermion description we fill the two sides of the inverted harmonic oscillator potential \[3,,4\].
To analyze this problem it is useful to realize that the asymptotic region of the weak coupling end in the target space geometry is associated to the asymptotic region of the Fermi sea far away from the maximum of the potential. So the two RR fluxes $`\nu ,\stackrel{~}{\nu }`$ that we discussed in section 2 are associated with the Fermi levels of the fermions on the two sides of the potential. Far from the maximum of the potential we can approximate the fermions as relativistic fermions since the depth of the Fermi sea is much larger than any finite energy we consider. It is also possible, and useful, to consider a basis for the inverted harmonic oscillator problem where the fermions are exactly relativistic . We review and extend this formalism in detail in Appendix A. There are actually two possible bases, which are naturally associated to the coordinates $`u=\frac{1}{\sqrt{2}}(px)`$ and $`s=\frac{1}{\sqrt{2}}(p+x)`$. These are the bases of in and out states and the S-matrix gives the relation between them. This relation is simply a Fourier transform.
So when we think about the asymptotic states we should think in terms of relativistic fermions. The asymptotic states live in the in and out Hilbert spaces of the fermions that are going towards the maximum of the potential or away from it. Each of these Hilbert spaces is described by two complex fermions
$$\psi _\pm ^{in},\psi ^{in\pm };\psi _\pm ^{out},\psi ^{out\pm }$$
The $`+/`$ indices denote fermions that are moving towards the right/left. It is very important not to confuse right and left moving matrix model fermions (denoted here by $`+/`$), which are moving to the right or left in eigenvalue space, with left and right movers in spacetime, which are related to incoming or outgoing states<sup>5</sup> Also, in the matrix model, do not confuse right and left moving fermions with fermions that are to the left and right side of the potential. For example, the in right moving fermion is to the left of the potential.. Our notation emphasizes the charge of the fermion under the $`U(1)`$ current which measures the number of right minus left moving fermions. This is the current associated to the scalar $`C`$ in spacetime. More precisely,
$$\begin{array}{cc}\hfill i(_t+_\varphi )C& \psi ^{in+}\psi _+^{in}\psi ^{in}\psi _{}^{in}\hfill \\ \hfill i(_t_\varphi )C& \psi ^{out+}\psi _+^{out}\psi ^{out}\psi _{}^{out}\hfill \end{array}$$
In principle we can specify freely the four Fermi levels of these four fermions. The fact that fermion number is conserved implies one relation between these four levels. So we have three independent levels which denote by $`\mu `$, $`\nu _{in}`$ and $`\nu _{out}`$. These are defined by saying that $`\mu \pm \nu _{in,out}`$ are the Fermi levels associated to the right and left moving incoming and outgoing fermions (see Figure 2).<sup>6</sup> When the time is rotated to Euclidean space we must also rotate $`\nu i\nu `$. This leads to an imaginary shift of the Fermi surface. This is consistent with the analysis of the $`\widehat{c}<1`$ systems which are similar to Euclidean $`\widehat{c}=1`$ where the RR flux was interpreted as an imaginary shift of the Fermi surface \[19--21\].
Fig. 2: Configurations with generic $`\nu _{in,out}`$ describe scattering amplitudes. Figures (a) and (b) describe scattering below the potential barrier $`\mu <0`$, and figures (c) and (d) describe scattering above the potential barrier $`\mu >0`$. Figures (a) and (c) describe the initial configurations, while (b) and (d) describe the final configuration. The dotted line represents the Fermi level characterized by $`\mu `$. (Even though we have represented the Fermi surface reaching all the way to the potential wall, we really should think of these configurations as asymptotic states, or as states in the in or out basis defined in the text.) Note that, as in the figure, the dominant scattering amplitudes for $`\mu >0`$ have $`\nu _{in}\nu _{out}`$; i.e. $`\stackrel{~}{\nu }0`$, while for $`\mu <0`$ they have $`\nu _{in}\nu _{out}`$; i.e. $`\nu 0`$. Note that in (b) $`\nu _{out}<0`$.
Fig. 3: In (a) we see an initial configuration of incoming fermions with $`\nu _{in}>0`$. In (b),(c) we see a outgoing configurations with $`\nu _{out}>0`$ and $`\nu _{out}<0`$ respectively. For the combination (a) (b) we have $`\nu _{out}=\nu _{in}`$ and therefore $`\stackrel{~}{\nu }=0`$, $`\nu =\nu _{in}`$. On the other hand for the combination (a), (c) we have $`\nu _{out}=\nu _{in}`$ or $`\nu =0`$, $`\stackrel{~}{\nu }=\nu _{in}`$.
In the case that we set $`|\nu _{in}||\nu _{out}|`$ we find that the incoming energy flux is not the same as the outgoing flux. Therefore we will need to add additional excitations. This is the same as in the discussion of the spacetime theory in section 2.
All the information about the inverted harmonic oscillator potential is contained in the map between in and out states (see Appendix A)
$$\psi _{a,ฯต}^{out}=\underset{b=\pm 1}{}๐ฎ_a^b(ฯต)\psi _{b,ฯต}^{in}=\frac{\mathrm{\Gamma }(\frac{1}{2}i\sqrt{2\alpha ^{}}ฯต)}{\sqrt{2\pi }}\underset{b=\pm 1}{}e^{i\frac{\pi }{2}ab(\frac{1}{2}i\sqrt{2\alpha ^{}}ฯต)}\psi _{b,ฯต}^{in}$$
where $`a,b=\pm 1`$ and $`\psi _{a,ฯต}^{in,out}`$ denote the annihilation operator for a fermion of energy $`ฯต`$.
We will now compute the transition amplitude between an in state with Fermi levels $`\mu `$, $`\nu _{in}`$ to an out state with Fermi levels $`\mu `$, $`\nu _{out}`$
$$๐=out(\mu ,\nu _{out})|in(\mu ,\nu _{in})$$
and interpret it as the partition function of the 0B theory with nonzero $`\nu `$ and $`\stackrel{~}{\nu }`$: $`๐=๐ต_{0B}(\mu ,\nu ,\stackrel{~}{\nu })`$.
For simplicity let us first assume that $`\nu _{in}=\nu _{out}=\nu >0`$. Consider the in state. The right moving fermions, which are created by $`\psi ^{in+}`$ are filled up to the Fermi level $`\mu +\nu `$, while the left moving fermions, created by $`\psi ^{in}`$ are filled up to the Fermi level $`\mu \nu `$. The same is true for the out fermions. So the overlap is given by<sup>7</sup> This expression for $`๐`$ suffers from a phase ambiguity. We can transform the incoming and outgoing, left and right moving Hilbert spaces by arbitrary energy independent phases. These four phases can be used to remove terms linear in $`\mu `$, $`\nu `$ (and later $`\stackrel{~}{\mu }`$) in $`\mathrm{log}๐`$.
$$๐=\underset{\mathrm{}<ฯต_n<\mu \nu }{}[๐ฎ_+^+(ฯต_n)๐ฎ_{}^{}(ฯต_n)๐ฎ_{}^+(ฯต_n)๐ฎ_+^{}(ฯต_n)]\underset{\mu \nu <ฯต_m<\mu +\nu }{}๐ฎ_+^+(ฯต_m)$$
where we have regularized the continuum by putting the system on a circle of length $`L`$ so that the density of states is $`dn=\frac{L}{2\pi }dฯต`$. We have used that up to the energy $`\mu \nu _{in}`$ both states are occupied. Note that the Fermi statistics produces the determinant of $`๐ฎ`$ for these states. On the other hand, for energies in the band between $`\mu \nu _{in}`$ and $`\mu +\nu _{in}`$ we have only the amplitude for an incoming right fermion going to an outgoing right fermion. Taking the logarithm of (4.1) and expressing the resulting sums in terms of integrals we obtain<sup>8</sup> We set $`\alpha ^{}=\frac{1}{2}`$.
$$\begin{array}{cc}\hfill \mathrm{log}๐=& \frac{L}{2\pi }[_\mathrm{\Lambda }^{\mu \nu }dฯต\mathrm{log}(\mathrm{\Gamma }(\frac{1}{2}iฯต)/\mathrm{\Gamma }(\frac{1}{2}+iฯต))+\hfill \\ & +_{\mu \nu }^{\mu +\nu }dฯต\mathrm{log}(\mathrm{\Gamma }(\frac{1}{2}iฯต)/\sqrt{2\pi })+\frac{\pi }{2}_{\mu \nu }^{\mu +\nu }dฯตฯต]\hfill \\ \hfill =& \frac{L}{2\pi }\left[_\mathrm{\Lambda }^{\mu +\nu }๐ฯต\mathrm{log}\left(\frac{\mathrm{\Gamma }(\frac{1}{2}iฯต)}{\sqrt{2\pi }}\right)_\mathrm{\Lambda }^{\mu \nu }๐ฯต\mathrm{log}\left(\frac{\mathrm{\Gamma }(\frac{1}{2}+iฯต)}{\sqrt{2\pi }}\right)+\pi \mu \nu \right]\hfill \end{array}$$
Here $`\mathrm{\Lambda }`$ is a cutoff on the bottom of the Fermi sea. Using
$$\mathrm{log}\left(\mathrm{\Gamma }(\frac{1}{2}+z)/\sqrt{2\pi }\right)=_0^{\mathrm{}}\frac{dt}{t}\left[\frac{e^{zt}}{2\mathrm{sinh}\frac{t}{2}}\frac{1}{t}+ze^t\right]$$
and defining $`\mathrm{\Xi }(y)`$ as related to the large radius limit of $`\mathrm{\Omega }(y,r)`$ of (3.1)
$$\mathrm{\Xi }(y)\underset{r\mathrm{}}{lim}\frac{\mathrm{\Omega }(y,r)}{2\pi r}=\frac{1}{2\pi }_0^{\mathrm{}}\frac{dt}{t}\left[e^{\frac{yt}{2}}\frac{1}{2t\mathrm{sinh}\frac{t}{2}}\frac{1}{t^2}+\frac{y}{2t}+(\frac{1}{24}\frac{y^2}{8})e^t\right]$$
(4.1) becomes
$$\mathrm{log}๐=iL\left[\mathrm{\Xi }(2i(\mu +|\nu |))\mathrm{\Xi }(2i(\mu |\nu |))+i\frac{1}{2}\mu |\nu |\right]$$
where we suppressed a $`\mathrm{\Lambda }`$ dependent imaginary constant and extended the answer to negative $`\nu `$ using the charge conjugation symmetry (left-right symmetry in the matrix model) $`\nu \nu `$.
Note that the 0B free energy at $`\nu =0`$ in the non-compact limit is given by
$$F_{0B}=\underset{\beta \mathrm{}}{lim}\frac{\mathrm{log}๐ต_{0B}}{\beta }=\underset{T_L\mathrm{}}{lim}\frac{\mathrm{log}๐}{iT_L}=\mathrm{\Xi }(2i\mu )\mathrm{\Xi }(2i\mu )$$
We have interpreted the cutoff $`L`$ as the length of Lorentzian time, $`T_L=L`$ since, for large $`L`$, this is the time it takes for the in state to come out from the scattering region. In other words the spatial cutoff $`L`$ is a good approximation for times which are of order $`L`$.
From the asymptotic expansion (3.1) we can find the asymptotic expansion of $`\mathrm{\Xi }`$.
$$\mathrm{\Xi }(y)=\frac{1}{2\pi }\left[(\mathrm{log}\frac{y}{2}\frac{3}{2})\frac{y^2}{8}\frac{1}{24}\mathrm{log}\frac{y}{2}\frac{7}{1440y^2}+๐ช(\frac{1}{y^4})\right]$$
This implies that the only perturbative terms in the imaginary part arise from the logarithms in (4.1)
$$\mathrm{Im}[\mathrm{\Xi }(2i\mu )]\frac{|\mu |\mu }{8}\left(1+\frac{1}{12}\frac{1}{\mu ^2}\right)$$
(Because of the dependence on the absolute value of $`\mu `$ they are not analytic and should be kept.) Thus, going back to (4.1), we find that the leading contribution to the real part of the log of the amplitude is
$$\begin{array}{cc}\hfill \mathrm{Re}[\mathrm{log}๐]& T_L[\frac{1}{8}|\mu +|\nu ||(\mu +|\nu |)(1+\frac{1}{12}\frac{1}{(\mu +|\nu |)^2})\hfill \\ & +\frac{1}{8}|\mu |\nu ||(\mu |\nu |)(1+\frac{1}{12}\frac{1}{(\mu |\nu |)^2})+\frac{1}{2}\mu |\nu |]\hfill \\ \hfill =& \{\begin{array}{cc}0\hfill & \mu \pm |\nu |>0\hfill \\ T_L\mu |\nu |\hfill & \mu \pm |\nu |<0\hfill \end{array}\hfill \end{array}$$
So we see that for $`\mu \pm |\nu |>0`$ the leading approximation to $`\mathrm{log}๐`$ is purely imaginary. Here we are in a configuration where the two Fermi seas are above the barrier (see figure 2 (c)/(d)), and the amplitudes is dominated by the leading order transmission over the barrier, as expected. On the other hand, for $`\mu \pm |\nu |<0`$, there is a negative real contribution. This implies that this processes is suppressed. Here the Fermi seas are below the barrier, and the in/out states are as in figure 3 (a)/(b). (Figure 2(a)/(b) depict a configuration with nonzero $`\stackrel{~}{\nu }`$.) In order to obey these boundary conditions we need to have tunneling processes. We need of the order of $`|\nu |`$ tunneling events per unit time. Each tunneling event contributes a factor of $`e^{\mu \pi }`$ (recall, $`\mu <0`$), which is the contribution of a charged D-instanton ZZ brane. The number of such factors depends on the total time $`T_L`$ as $`T_L|\nu |/\pi `$. This leads to (4.1). Note that the effects of the instantons do not exponentiate since we are looking at a very special process where only a definite number of instantons could contribute. Of course one could also study processes where $`\mu |\nu |<0<\mu +|\nu |`$. Then, the second term in (4.1) also contributes.
The expression (4.1) for the partition function $`๐`$ has a few interesting consequences. The two terms $`\mathrm{\Xi }(2i(\mu +|\nu |)`$ and $`\mathrm{\Xi }(2i(\mu |\nu |)`$ can be interpreted as the contribution of the fermions in the left and the right side of the potential. Therefore, either $`\mathrm{\Xi }(2i\mu )`$ or $`\mathrm{\Xi }(2i\mu )`$ can be viewed as a nonperturbative definition of the free energy of the bosonic $`c=1`$ system with Fermi level $`\mu `$. (Recall, the bosonic system has fermions only in one side of the potential.) From this perspective the problem with the $`c=1`$ system is that $`\mathrm{\Xi }(\pm 2i\mu )`$ are complex. The sign ambiguity in the definition changes the sign of the imaginary part which signals the instability of the system.
The expression (4.1) also gives an intuitive explanation of the holomorphic factorization we have seen before. Up to the simple term which depends on $`\mu |\nu |`$ the partition function factorizes as a product of $`\mathrm{exp}(iL\mathrm{\Xi }(2i(\mu +|\nu |)`$ and $`\mathrm{exp}(iL\mathrm{\Xi }(2i(\mu |\nu |)`$ which are associated with the incoming fermions from the left and the right side of the potential. Our definition of the RR-flux is such that it does not mix these two kids of fermions, and therefore we can specify independent Fermi levels for them, $`\mu \pm |\nu |`$. In Euclidean space $`|\nu |i|\nu |`$ and therefore this separation explains the holomorphic factorization we discussed above. We will soon add nonzero $`\stackrel{~}{\nu }`$, and will study the problem with a Euclidean time circle. The separation of these modes will persist. It underlies the holomorphic factorization of the partition function.
Repeating this analysis for $`\nu _{in}=\nu _{out}=\stackrel{~}{\nu }`$ we find an answer that is very similar to (4.1) except that the last term changes sign
$$\mathrm{log}๐=iL\left[\mathrm{\Xi }(2i(\mu +|\stackrel{~}{\nu }|))\mathrm{\Xi }(2i(\mu |\stackrel{~}{\nu }|))i\frac{1}{2}\mu |\stackrel{~}{\nu }|\right]$$
In this case the real part is small for sufficiently large negative $`\mu `$ but it is behaves as $`L\mu |\stackrel{~}{\nu }|`$ for large positive $`\mu `$. This is consistent with the duality symmetry $`\mu \mu `$, $`\nu \stackrel{~}{\nu }`$.
In the case that $`\nu _{in}^2\nu _{out}^2`$ we have to insert extra asymptotic states in order to balance the energy flux. We will do this in more detail in the Euclidean computation in the next subsection.
4.2. Computation at finite $`R`$
In this section we consider the finite temperature partition function, where the time direction has period $`\beta =2\pi R`$. Configurations with RR fluxes correspond to configurations where the field $`C`$ or its dual have winding along the Euclidean time direction. We have said that the field $`C`$ is at the self dual radius, $`CC+2\pi \sqrt{2}`$ with the normalization (2.1). This implies the following quantization condition for the constant part of the Euclidean time derivative
$$_\tau C=i2\sqrt{2}\stackrel{~}{\nu }=\sqrt{2}\frac{\stackrel{~}{q}}{R},\stackrel{~}{q}๐$$
We can similarly think about the quantization condition for, $`\stackrel{~}{C}`$, the dual of $`C`$. This gives
$$_\varphi C=2\sqrt{2}\nu =i\sqrt{2}\frac{q}{R},q๐$$
We then define
$$q_{in}=q+\stackrel{~}{q},q_{out}=q\stackrel{~}{q}$$
Note that $`q,\stackrel{~}{q},q_{in},q_{out}`$ are integer, but $`q_{in}q_{out}`$ is always even. In these conventions an in right moving fermion has $`q_{in}=1`$. Note that we can consider in states where we have a left moving spin field and a right moving spin field. If the charge of the spin fields are opposite this configuration gives us $`q_{in}=1`$. But then we should also have spin fields in the out state since (4.1) (4.1) (4.1) imply that $`q^{in,out}`$ are either both odd or both even.
There are two closely related ways of thinking about the Euclidean computation. One is to view it as an analytic continuation of the Lorentzian scattering computation (4.1). The only difference is that the asymptotic regions now look like a cylinder. So we think of the in and out states as living on a cylinder and we expand them in fourier modes along the compact direction. The analytic continuation of (4.1) gives the relation between in and out fields. The resulting relation can be summarized as (we continue to set $`\alpha ^{}=\frac{1}{2}`$)
$$\begin{array}{cc}& \psi _r^{outa}\psi _{b,s}^{in}=\delta _{r,s}\frac{\mathrm{\Gamma }(\frac{1}{2}+i\mu +s/R)}{\sqrt{2\pi }}e^{\frac{\pi }{2}(\mu i\frac{s}{R})abi\frac{\pi }{4}ab}\hfill \\ & \psi _{a,r}^{out}\psi _s^{inb}=\delta _{r,s}\frac{\mathrm{\Gamma }(\frac{1}{2}i\mu +s/R)}{\sqrt{2\pi }}e^{\frac{\pi }{2}(\mu +i\frac{s}{R})ab+i\frac{\pi }{4}ab}\hfill \end{array}$$
where $`r,s๐+\frac{1}{2}`$, $`r,s>0`$. When we do this analytic continuation of (4.1) we might be a bit unsure about the sign for $`s`$ in the right hand side. This sign is determined by doing the analytic continuation of the fields carefully and demanding that the mode, $`\psi _s^{in}`$ does not annihilate the vacuum, $`\psi _s|00`$ for positive $`s`$ (in our conventions $`\psi _s^{in/out}|0=0`$ for $`s>0`$). In this description we are interested in computing an inner product of the form
$$\mathrm{\Psi }_{out}|\mathrm{\Psi }_{in}$$
Notice that from the target space viewpoint, the states $`\mathrm{\Psi }_{in}`$ and $`\mathrm{\Psi }_{out}`$ are determined by choosing the non-normalizable behavior of the anti-holomporphic and the holomorphic parts of the target space fields $`T`$ and $`C`$ near the boundary. As usual, this correspondence involves bosonization of the fermions and an identification of the modes of the bosonic field with the modes of the $`T`$ and $`C`$ fields.
The other way of thinking about the problem consists in viewing the problem as defined on a half cylinder where the in and out fields are antiholomorphic and holomorphic fields respectively. We impose boundary condition at the asymptotic end of the cylinder by specifying a state in the Hilbert space for fermions on a cylinder. In the capped end of the semi infinite cylinder we insert a boundary state which encodes the effects of the scattering amplitude. This boundary state is also computed by analytically continuing (4.1). It is clear that both pictures are equivalent and which one we choose is a matter of taste.
Let us consider a configuration with general $`q_{in}`$ and $`q_{out}`$, and first consider the case
$$q_{in}q_{out}0$$
For simplicity, let us limit ourselves to $`q_{in},q_{out}2๐`$. <sup>9</sup> This case is simpler because we do not need to introduce spin fields. So we will have a state in the in Hilbert space of dimension $`\mathrm{\Delta }_{in}`$. Since our problem is invariant under translations in Euclidean time, we need that $`\mathrm{\Delta }_{out}=\mathrm{\Delta }_{in}`$. We will be interested in considering the state with lowest $`\mathrm{\Delta }_{in}`$ since this is the state that corresponds to exciting only the constant part of the RR field strength (the gradient of $`(_\varphi i_t)Cq_{in}`$). This lowest dimension state is
$$\mathrm{\Delta }_{in}=\frac{q_{in}^2}{4}=\mathrm{\Delta }_{out}$$
which corresponds to the state
$$|\mathrm{\Psi }_{in}=\underset{l=1}{\overset{q_{in}/2}{}}\psi _{l+1/2}^{in+}\psi _{,l+\frac{1}{2}}^{in}|0$$
States with higher dimension, with the same $`q_{in}`$ correspond to exciting other oscillator modes of $`T`$ and $`C`$.
In the out Hilbert space we need a state with the same conformal dimension. Since the S-matrix is given by the product of one body S-matrices it is clear that we need as many fermions and holes in the out Hilbert space as we have in the in Hilbert space. Since our system is time translation invariant, the amplitudes are diagonal in the mode number. So the non-zero amplitudes have the form
$$0|\underset{l=1}{\overset{q_{in}/2}{}}\psi _{l\frac{1}{2}}^{outa_n}\psi _{b_n,l\frac{1}{2}}^{out}\psi _{l+\frac{1}{2}}^{in+}\psi _{,l+\frac{1}{2}}^{in}|0$$
Since the charge of the out state has to be $`q_{out}`$ we need that $`b_na_n=q_{out}`$<sup>10</sup> Notice that we are defining the charge as the charge of the ket, which is minus the charge of the out operators explicitly appearing in (4.1).. There are several ways to assign values of $`a_n`$ and $`b_n`$. Of the $`q_{in}`$ out-fermions, $`q=(q_{in}+q_{out})/2`$ should have $`out`$ charge minus one and $`\stackrel{~}{q}=(q_{in}q_{out})/2`$ should have $`out`$ charge plus one. The number of possibilities of achieving this is
$$N(q,\stackrel{~}{q})=\frac{q_{in}!}{[\frac{1}{2}(q_{in}+q_{out})]![\frac{1}{2}(q_{in}q_{out})]!}=\frac{(q+\stackrel{~}{q})!}{q!\stackrel{~}{q}!}$$
Using (4.1) we can compute (4.1) and obtain
$$\begin{array}{cc}\hfill ๐(\mu ,q_{in},q_{out})& =e^{i\phi (a_n,b_n,R)}\underset{n=1}{\overset{q_{in}/2}{}}\frac{|\mathrm{\Gamma }(\frac{1}{2}i\mu +(n\frac{1}{2})/R)|^2}{2\pi }e^{\frac{\pi \mu (a_nb_n)}{2}}๐(\mu ,0,0)\hfill \\ & =e^{i\phi (a_n,b_n,R)}e^{\frac{\pi \mu q_{out}}{2}}\underset{n=1}{\overset{q_{in}/2}{}}\frac{|\mathrm{\Gamma }(\frac{1}{2}i\mu +(n\frac{1}{2})/R)|^2}{2\pi }๐(\mu ,0,0)\hfill \end{array}$$
where we have used $`q_{out}=(b_na_n)`$. Note that up to the phase $`e^{i\phi (a_n,b_n,R)}`$ the answer does not depend on the particular operator among all the $`N(q_{in},q_{out})`$ operators in (4.1) with the same charges<sup>11</sup> The technical reason for this is the fact that the difference between the right to right vs right to left amplitudes (4.1) is a simple exponential.. It will be important below that this phase is independent of $`\mu `$.
It is straightforward to extend the computation (4.1) to values of $`q_{in}`$ and $`q_{out}`$ which do not satisfy (4.1). The answer is expressed most easily in terms of $`q`$ and $`\stackrel{~}{q}`$
$$๐(\mu ,q_{in},q_{out})=e^{i\phi (a_n,b_n,R)}e^{\frac{\pi \mu }{2}(|q||\stackrel{~}{q}|)}\underset{n=1}{\overset{\widehat{q}/2}{}}\frac{|\mathrm{\Gamma }(\frac{1}{2}+i\mu +\frac{n\frac{1}{2}}{R})|^2}{2\pi }๐ต_{0B}(\mu ,q=\stackrel{~}{q}=0,R)$$
where $`\widehat{q}=|q|+|\stackrel{~}{q}|=q_{in}`$ for the case (4.1).
Using (4.1) and the expression for $`\mathrm{\Omega }(y,r)`$ (3.1) we find for even $`\widehat{q}=2k`$ <sup>12</sup> This is the same as the recursion relations found for the function $`G(x)`$ in .
$$\begin{array}{cc}\hfill \mathrm{\Omega }(y=\frac{2k}{R}+2i\mu ,R)& \mathrm{\Omega }(y=0+2i\mu ,R)=\hfill \\ & =_0^{\mathrm{}}\frac{dt}{t}\left[e^{i\mu t}\frac{e^{\frac{kt}{R}}1}{4\mathrm{sinh}\frac{t}{2}\mathrm{sinh}\frac{t}{2R}}+\frac{k}{t}(\frac{k^2}{2R}+ik\mu )e^t\right]\hfill \\ & =\underset{n=1}{\overset{k}{}}_0^{\mathrm{}}\frac{dt}{t}\left[\frac{e^{(\frac{2n1}{2R}+i\mu )t}}{2\mathrm{sinh}\frac{t}{2}}\frac{1}{t}+(\frac{2n1}{2R}+i\mu )e^t\right]\hfill \\ & =\underset{n=1}{\overset{k}{}}\mathrm{log}\left(\frac{\mathrm{\Gamma }(\frac{1}{2}+i\mu +\frac{n\frac{1}{2}}{R})}{\sqrt{2\pi }}\right)\hfill \end{array}$$
Using this relation and the expression for $`๐ต_{0B}(\mu ,q=\stackrel{~}{q}=0)`$ (see Appendix A), (4.1) can be written as
$$\begin{array}{cc}\hfill \mathrm{log}๐ต_{0B}(\mu ,q,\stackrel{~}{q},R)=& \mathrm{log}๐(\mu ,q_{in},q_{out})=i\phi (a_n,b_n,R)+\frac{\pi \mu }{2}(|q||\stackrel{~}{q}|)+\hfill \\ & +\mathrm{\Omega }(y=\frac{\widehat{q}}{R}+2i\mu ,R)+\mathrm{\Omega }(y=\frac{\widehat{q}}{R}2i\mu ,R)\hfill \end{array}$$
which is our final expression for the free energy.
After we apply T-duality
$$R_B=\frac{\alpha ^{}}{R_A},\mu _B=\frac{R_A}{\sqrt{2\alpha ^{}}}\mu _A$$
we find that (4.1) becomes the same as (3.1), up to the phase and analytic terms in $`\mu `$ proportional to $`\mathrm{log}R`$. These terms are related to the fact that we need to change the UV cutoff $`\mathrm{\Lambda }`$ when we perform T-duality (see the appendix).
Note that this 0B computation produces naturally the term that involves $`(|q||\stackrel{~}{q}|)`$ while in the 0A problem we had to introduce this term โby handโ in order to match the expected asymptotic behavior.
Notice that this procedure produces answers which are consistent with T-duality, while previous studies did not.
The answer (4.1) has the expected symmetries: $`A(\mu ,q,\stackrel{~}{q})=A(\mu ,q,\stackrel{~}{q})=A(\mu ,\stackrel{~}{q},q)=A(\mu ,q,\stackrel{~}{q})^{}`$. They follow from the two $`๐_2`$ symmetries and time reversal.
Acknowledgements
It is a pleasure to thank I. Klebanov and E. Witten for discussions. This work was supported in part by grant #DE-FG02-90ER40542.
Appendix A. Chiral Quantization of Matrix Models
The purpose of this appendix is to derive some known results about the $`\widehat{c}=1`$ 0A and 0B matrix models. We will use a formalism which was first introduced in and was later used and elaborated on in \[39--47\]. It highlights the chiral nature of the problem and the scattering from and to null infinities. One of the advantages of this formalism is that the theory is expressed in terms of free relativistic fermions. The nontrivial scattering appears as a nonlocal transform between the incoming and the outgoing descriptions. The parabolic cylinder functions of the inverted harmonic oscillator are replaced by simple wave functions in the $`p\pm x`$ representation. $`p\pm x`$ are the analog of creation and annihilation operators of the ordinary harmonic oscillator and their eigenstates are analogous to the familiar coherent states. However, unlike the ordinary harmonic oscillator, since $`p\pm x`$ are hermitian operators, their eigenvalues are real and the inner product of functions in these representations is standard.
We will present this formalism, will clarify some of its properties and will extend it. We will start the discussion of the first quantized theories with some general properties of eigenstates of $`p\pm x`$, and will then use them in the special cases relevant to the 0B and 0A strings. Then, we will study the second quantized theories and will compute their free energies.
A.1. First quantized problems
A single upside down harmonic oscillator
Consider first a generic quantum mechanical problem of a single degree of freedom. Standard bases of orthonormal states are $`|x`$ and $`|p`$, which are coordinate and momentum eigenstates respectively, with $`x|p=\frac{1}{\sqrt{2\pi }}e^{ipx}`$. We will also be interested in the bases $`|s`$ and $`|u`$ which are orthonormal eigenstates of
$$\begin{array}{cc}& S=\frac{P+X}{\sqrt{2}}\hfill \\ & U=\frac{PX}{\sqrt{2}}\hfill \\ & [S,U]=i\hfill \end{array}$$
Fig. 4: The phase space for a harmonic oscillator parameterized by the $`x,p`$ coordinates or the $`s,u`$ coordinates. The dotted lines denote various possible classical trajectories.
It is easy to find the inner products
$$\begin{array}{cc}& x|s=\frac{2^{\frac{1}{4}}e^{i\frac{\pi }{8}}}{\sqrt{2\pi }}\mathrm{exp}\left(i(\frac{x^2}{2}+\sqrt{2}sx\frac{s^2}{2})\right)\hfill \\ & x|u=\frac{2^{\frac{1}{4}}e^{i\frac{\pi }{8}}}{\sqrt{2\pi }}\mathrm{exp}\left(i(\frac{x^2}{2}+\sqrt{2}ux+\frac{u^2}{2})\right)\hfill \\ & s|u=\frac{1}{\sqrt{2\pi }}\mathrm{exp}\left(isu\right)\hfill \end{array}$$
The $`s`$ and $`u`$ dependent phases in $`|s`$ and $`|u`$ are such that $`U=S\sqrt{2}X`$ acts on $`s|u`$, and $`s|x`$ as $`i_s`$ and similarly for the action of $`S`$ on $`u|s`$ and $`u|x`$. In the last expression we defined the integral $`๐xe^{ix^2}`$ as $`๐xe^{(i0^+)x^2}=\sqrt{i\pi }`$. We have chosen the constant phases of the first two lines so as to simplify the last line and some of the subsequent formulas.
So, let us now focus on the inverted harmonic oscillator with the Lagrangian and Hamiltonian
$$\begin{array}{cc}& L=\frac{1}{2}(\dot{X}^2+X^2)\hfill \\ & H=\frac{1}{2}(P^2X^2)=\frac{1}{2}(SU+US)\hfill \end{array}$$
It is easy to work out the time evolution
$$\begin{array}{cc}& e^{iHt}|s=e^{\frac{t}{2}}|e^ts\hfill \\ & e^{iHt}|u=e^{\frac{t}{2}}|e^tu\hfill \\ & s|e^{iHt}=e^{\frac{t}{2}}e^ts|\hfill \\ & u|e^{iHt}=e^{\frac{t}{2}}e^tu|\hfill \\ & s|e^{iHt}|u=\frac{1}{\sqrt{2\pi }}e^{\frac{t}{2}}\mathrm{exp}\left(isue^t\right)\hfill \end{array}$$
Here the factors of $`e^{\pm \frac{t}{2}}`$ are needed for unitarity, but also come out of (A.1) by writing, in the $`s`$ basis $`H=is_s\frac{i}{2}`$, and $`H=iu_u+\frac{i}{2}`$ in the $`u`$ basis.
The operators (A.1) are similar to the creation and annihilation operators of the ordinary (or โupside upโ) harmonic oscillator. One difference is that in our case, $`S`$ and $`U`$ are hermitian operators and not hermitian conjugates to each other. Correspondingly the states $`|s`$ or $`|u`$ are analogous to coherent states. Since these operators are hermitian the states $`s|`$ and $`u|`$ will also eigenstates of $`S`$ and $`U`$ respectively. These two bases will be useful to describe the initial and final states of the upside down harmonic oscillator. In other words, the incoming states will be naturally described in terms of the $`u`$ basis and the outgoing states in terms of the $`s`$ basis. This can be seen quite naturally by looking at the shape of trajectories in fig. 4, but will be seen more precisely later.
There are two linearly independent energy eigenstates for every energy $`ฯต`$. In the $`x`$ representation the wavefunctions are the two parabolic cylinder functions. They can be taken to be even and odd under the parity transformation $`xx`$. Alternatively, we can take one wavefunction to correspond to a wave coming from the left and scattered to the right and back to the left, and the other wave function obtained from this one by $`xx`$.
In the $`s`$ representation the energy eigenstates with eigenvalue $`ฯต`$ are $`s^{iฯต\frac{1}{2}}`$. The singularity at $`s=0`$ leads to a two fold doubling of the number of states $`|ฯต,out\pm `$, where the label $`out`$ will be explained shortly. Their wavefunctions are
$$\begin{array}{cc}& s|ฯต,out+=\{\begin{array}{cc}\frac{1}{\sqrt{2\pi }}s^{iฯต\frac{1}{2}}\hfill & s>0\hfill \\ 0\hfill & s<0\hfill \end{array}\hfill \\ & s|ฯต,out=\{\begin{array}{cc}0\hfill & s>0\hfill \\ \frac{1}{\sqrt{2\pi }}(s)^{iฯต\frac{1}{2}}\hfill & s<0\hfill \end{array}\hfill \end{array}$$
By looking at the trajectories in fig. 4 we see that $`|ฯต,out+`$ states are states that in their outgoing modes contain only a right moving piece. While the states $`|ฯต,out`$ contain only a left moving piece in their outgoing modes. Therefore we will refer to them as โout statesโ.
Another natural basis arises from the $`u`$ representation
$$\begin{array}{cc}& u|ฯต,in+=\{\begin{array}{cc}\frac{1}{\sqrt{2\pi }}u^{iฯต\frac{1}{2}}\hfill & u>0\hfill \\ 0\hfill & u<0\hfill \end{array}\hfill \\ & u|ฯต,in=\{\begin{array}{cc}0\hfill & u>0\hfill \\ \frac{1}{\sqrt{2\pi }}(u)^{iฯต\frac{1}{2}}\hfill & u<0\hfill \end{array}\hfill \end{array}$$
These are states which contain only right/left moving incoming pieces for $`+/`$. We will refer to them as โin statesโ.
Semiclassically the incoming states have $`u\pm \mathrm{}`$ and $`s0`$, while the outgoing states have $`u0`$ and $`s\pm \mathrm{}`$. Therefore, it is natural to take the incoming states to be $`|ฯต,in\pm `$, where $`|ฯต,in+`$ describes a particle coming from the left (negative $`x`$) and $`|ฯต,in`$ describes a particle coming from the right (positive $`x`$). Similarly, the outgoing states are $`|ฯต,out\pm `$. Here, $`|ฯต,out+`$ describes a particle going to the right (positive $`x`$), and $`|ฯต,out`$ describes a particle going to the left (negative $`x`$).
These two bases are related by a unitary transformation
$$\begin{array}{cc}& \left(\begin{array}{c}|ฯต,out+\\ |ฯต,out\end{array}\right)=๐ฎ\left(\begin{array}{c}|ฯต,in+\\ |ฯต,in\end{array}\right)\hfill \\ & ๐ฎ=\left(\begin{array}{cc}\frac{e^{i\frac{\pi }{4}}e^{\frac{\pi ฯต}{2}}}{\sqrt{2\pi }}\mathrm{\Gamma }(\frac{1}{2}iฯต)& \frac{e^{i\frac{\pi }{4}}e^{\frac{\pi ฯต}{2}}}{\sqrt{2\pi }}\mathrm{\Gamma }(\frac{1}{2}iฯต)\\ \frac{e^{i\frac{\pi }{4}}e^{\frac{\pi ฯต}{2}}}{\sqrt{2\pi }}\mathrm{\Gamma }(\frac{1}{2}iฯต)& \frac{e^{i\frac{\pi }{4}}e^{\frac{\pi ฯต}{2}}}{\sqrt{2\pi }}\mathrm{\Gamma }(\frac{1}{2}iฯต)\end{array}\right)=\left(\begin{array}{cc}\frac{e^{i\frac{\pi }{4}}e^{i\mathrm{\Phi }_B(ฯต)}}{\sqrt{1+e^{2\pi ฯต}}}& \frac{e^{i\frac{\pi }{4}}e^{i\mathrm{\Phi }_B(ฯต)}}{\sqrt{1+e^{2\pi ฯต}}}\\ \frac{e^{i\frac{\pi }{4}}e^{i\mathrm{\Phi }_B(ฯต)}}{\sqrt{1+e^{2\pi ฯต}}}& \frac{e^{i\frac{\pi }{4}}e^{i\mathrm{\Phi }_B(ฯต)}}{\sqrt{1+e^{2\pi ฯต}}}\end{array}\right)\hfill \\ & e^{i\mathrm{\Phi }_B(ฯต)}=\sqrt{\frac{\mathrm{\Gamma }(\frac{1}{2}iฯต)}{\mathrm{\Gamma }(\frac{1}{2}+iฯต)}}\hfill \end{array}$$
Here we wrote $`ฯต,out+|ฯต,in=๐s๐uฯต,out+|ss|uu|ฯต,in`$ and we used (A.1), (A.1), (A.1). Another way to understand these bases is to express the $`out`$ states as functions of $`u`$ and the $`in`$ states as functions of $`s`$. Using (A.1), or more directly by using $`s|u`$ in (A.1) and Fourier transforming (A.1)(A.1), we find
$$\begin{array}{cc}& s|ฯต,in+=\frac{e^{i\frac{\pi }{4}}e^{i\mathrm{\Phi }_B(ฯต)}}{\sqrt{2\pi }\sqrt{1+e^{2\pi ฯต}}}(s+i0^+)^{iฯต\frac{1}{2}}\hfill \\ & s|ฯต,in=\frac{e^{i\frac{\pi }{4}}e^{i\mathrm{\Phi }_B(ฯต)}}{\sqrt{2\pi }\sqrt{1+e^{2\pi ฯต}}}(s+i0^{})^{iฯต\frac{1}{2}}\hfill \\ & u|ฯต,out+=\frac{e^{i\frac{\pi }{4}}e^{i\mathrm{\Phi }_B(ฯต)}}{\sqrt{2\pi }\sqrt{1+e^{2\pi ฯต}}}(u+i0^{})^{iฯต\frac{1}{2}}\hfill \\ & u|ฯต,out=\frac{e^{i\frac{\pi }{4}}e^{i\mathrm{\Phi }_B(ฯต)}}{\sqrt{2\pi }\sqrt{1+e^{2\pi ฯต}}}(u+i0^+)^{iฯต\frac{1}{2}}\hfill \end{array}$$
where the $`s+i0^\pm `$ prescription means that for negative $`s`$ we substitute $`s^{i\alpha }=(|s|e^{\pm i\pi })^{i\alpha }=(s)^{i\alpha }e^{\pi \alpha }`$. So the $`|ฯต,in\pm `$ states in the $`s`$ representation $`s|ฯต,in\pm `$ are linear combinations of $`s|ฯต,out\pm `$ whose precise coefficients are determined by thinking of $`s|ฯต,in\pm `$ as functions that are analytic in the upper/lower half $`s`$ plane. The situation is very similar to one that arises in the physics of Rindler space when we express the Minkowski wavefunctions in terms of the Rindler wavefunctions. In our case this arises from the fact that $`|ฯต,in+`$ has support only for $`u>0`$, this implies that in the $`s`$ representation this is an analytic function for $`Im(s)>0`$. In the Rindler case, the positive frequency condition on the Minkowski wavefunctions implies similar analyticity properties. In a connection between these fermions and fermions in de-Sitter space was studied. In the thermal looking nature of these amplitudes was explored.
The even and odd states $`|ฯต,out+\pm |ฯต,out`$ and $`|ฯต,in+\pm |ฯต,in`$ diagonalize (A.1)
$$\begin{array}{cc}& |ฯต,in+\pm |ฯต,in=e^{i\phi _\pm (ฯต)}\left(|ฯต,out+\pm |ฯต,out\right)\hfill \\ & e^{i\phi _+(ฯต)}=e^{i\mathrm{\Phi }_B(ฯต)}\frac{e^{i\frac{\pi }{4}}+e^{i\frac{\pi }{4}}e^{\pi ฯต}}{\sqrt{1+e^{2\pi ฯต}}}=2^{iฯต}\frac{\mathrm{\Gamma }(\frac{1}{4}\frac{iฯต}{2})}{\mathrm{\Gamma }(\frac{1}{4}+\frac{iฯต}{2})}\hfill \\ & e^{i\phi _{}(ฯต)}=e^{i\mathrm{\Phi }_B(ฯต)}\frac{e^{i\frac{\pi }{4}}e^{i\frac{\pi }{4}}e^{\pi ฯต}}{\sqrt{1+e^{2\pi ฯต}}}=2^{iฯต}i\frac{\mathrm{\Gamma }(\frac{3}{4}\frac{iฯต}{2})}{\mathrm{\Gamma }(\frac{3}{4}+\frac{iฯต}{2})}\hfill \end{array}$$
Our interpretation of (A.1)(A.1) as in and out states is further supported by comparing (A.1)(A.1) and (A.1). The in states $`|ฯต,in\pm `$ have support only for one sign of $`u`$ and for both signs of $`s`$. This is the expected behavior of incoming states. The $`out`$ states $`|ฯต,out\pm `$ have support only for one sign of $`s`$ and for both signs of $`u`$. This is the expected behavior of outgoing states. Furthermore, for large $`|ฯต|`$ the relation between the two bases is simple. Up to an $`ฯต`$ dependent phase $`|ฯต,in\pm |ฯต,out\pm `$ for positive $`ฯต`$ and $`|ฯต,in\pm |ฯต,out`$ for negative $`ฯต`$. This is consistent with the semiclassical picture of complete transmission for positive $`ฯต`$ and complete reflection for negative $`ฯต`$.
One can actually show more precisely why it is reasonable associate the basis $`|ฯต,in\pm `$ with incoming states. For that purpose we can compute $`x|ฯต,in\pm `$ using (A.1) (A.1) . We do not need the exact answer, which is a combination of parabolic cylinder functions. We only need the behavior of the function for large $`x`$. Since $`x`$ is large we can compute the answer by saddle point integration. The saddle point equation for $`u`$ is $`\sqrt{2}x+uฯต/u0`$. The two saddle points, at $`u\sqrt{2}x`$ and at $`uฯต/(\sqrt{2}x)`$, give the incoming and outgoing pieces of the $`x`$ space wavefunction, which go like $`e^{i\frac{x^2}{2}}`$ and $`e^{i\frac{x^2}{2}}`$ to leading order in $`x`$. Note that the first saddle point arises only for $`\pm x<0`$ for $`|ฯต,in\pm `$. This means that the incoming wavefunction is supported to the left/right side of the potential for $`|ฯต,in\pm `$. Furthermore, the coefficient of the first saddle point is energy independent (except for a simple, expected, factor of $`|x|^{iฯต\frac{1}{2}}`$). This is the natural normalization for the incoming states. The integral in the region close to the second saddle point gives us the reflected part of the wavefunction and contains the information about scattering phase.
Repeating this discussion for $`|ฯต,out\pm `$ we can understand why it is natural to associate them to outgoing states which are right or left moving.
Two upside down harmonic oscillators
Now, let us discuss the same problem but with two degrees of freedom $`X_1`$ and $`X_2`$ and their conjugate momenta $`P_1`$ and $`P_2`$. We change variables to polar coordinates $`X_1=X\mathrm{cos}\theta `$, $`X_2=X\mathrm{sin}\theta `$. The momentum conjugate to $`\theta `$ is
$$q=X_1P_2X_2P_1$$
and we can work in a sector where it is a fixed integer $`c`$ number. Then, we have two natural bases of states $`|x`$ and $`|p`$ which are eigenstates of $`X`$ and $`P=\mathrm{cos}\theta P_1+\mathrm{sin}\theta P_2`$ respectively. Note that, even though $`(X_i,P_i)`$ are canonically conjugate, $`P`$ is not the momentum conjugate to $`X`$.
As in (A.1), we define
$$\begin{array}{cc}& S_i=\frac{P_i+X_i}{\sqrt{2}}\hfill \\ & U_i=\frac{P_iX_i}{\sqrt{2}}\hfill \\ & [S_i,U_j]=i\delta _{i,j}\hfill \end{array}$$
and we can again change to โpolar coordinatesโ
$$\begin{array}{cc}& S_1=S\mathrm{cos}\theta _s\hfill \\ & S_2=S\mathrm{sin}\theta _s\hfill \\ & U_1=U\mathrm{cos}\theta _u\hfill \\ & U_2=U\mathrm{sin}\theta _u\hfill \end{array}$$
The momenta conjugate to $`S`$ and $`U`$ are $`P_s`$ and $`P_u`$. It is important that they are not given by $`U`$ and $`S`$. However, $`q`$ of (A.1) can be written also as $`q=S_1U_2S_2U_1`$, and it is the momentum conjugate to both $`\theta `$, $`\theta _s`$ and $`\theta _u`$. This follows from the fact that this is the charge of the same rotation symmetry.
Let us study the various bases in more detail. The simplest states are $`X_i`$ eigenstates, $`|x_1,x_2`$, or in polar coordinates $`|x,\theta `$. Note that the eigenvalue $`x`$ is positive. They satisfy
$$x_1,x_2|x,\theta =\sqrt{x}\delta (x_1x\mathrm{cos}\theta )\delta (x_2x\mathrm{sin}\theta )$$
Instead of diagonalizing $`\theta `$ it is better to diagonalize $`q`$, $`|x,q=\frac{1}{\sqrt{2\pi }}๐\theta e^{iq\theta }|x,\theta `$. We will use similar notation for various bases diagonalizing the $`S`$ or $`U`$ variables (again, the eigenvalues $`s`$ and $`u`$ are positive). One way to find the inner products between these bases is to convert to the bases where the Cartesian coordinates $`X_i`$, $`S_i`$ or $`U_i`$ are diagonal and then use (A.1) for each of them. We readily find (recall, $`_0^{2\pi }๐\theta e^{iq\theta }e^{ia\mathrm{cos}\theta }=2\pi i^qJ_q(a)`$)
$$\begin{array}{cc}\hfill s,\theta _s|x,q& =\frac{e^{i\frac{\pi }{4}}i^{\frac{3q}{2}}\sqrt{2xs}}{(2\pi )^{\frac{3}{2}}}_0^{2\pi }๐\theta e^{iq\theta }\mathrm{exp}\left(i(\frac{x^2}{2}\sqrt{2}sx\mathrm{cos}(\theta \theta _s)+\frac{s^2}{2})\right)\hfill \\ & =\frac{e^{i\frac{\pi }{4}}(1)^qi^{\frac{q}{2}}\sqrt{2xs}}{\sqrt{2\pi }}\mathrm{exp}\left(iq\theta _s+i\frac{x^2}{2}+i\frac{s^2}{2}\right)J_q(\sqrt{2}sx)\hfill \\ \hfill u,\theta _u|x,q& =\frac{e^{i\frac{\pi }{4}}i^{\frac{q}{2}}\sqrt{2xu}}{(2\pi )^{\frac{3}{2}}}_0^{2\pi }๐\theta e^{iq\theta }\mathrm{exp}\left(i(\frac{x^2}{2}\sqrt{2}ux\mathrm{cos}(\theta \theta _u)\frac{u^2}{2})\right)\hfill \\ & =\frac{e^{i\frac{\pi }{4}}(1)^qi^{\frac{q}{2}}\sqrt{2xu}}{\sqrt{2\pi }}\mathrm{exp}\left(iq\theta _ui\frac{x^2}{2}i\frac{u^2}{2}\right)J_q(\sqrt{2}ux)\hfill \\ \hfill s,q|u,q^{}& =\frac{i^q\delta _{q,q^{}}\sqrt{su}}{2\pi }_0^{2\pi }๐\theta _ue^{iq\theta _u+isu\mathrm{cos}\theta _u}=\sqrt{su}J_q(su)\delta _{q,q^{}}\hfill \end{array}$$
where we have chosen the overall phases to simplify some formulas later. The first two expressions demonstrate that $`q`$ is the momentum conjugate to $`\theta `$, $`\theta _s`$ and $`\theta _u`$. The last inner product can be derived in several different ways whose consistency relies on (or better, gives a proof of) the Weberโs formula
$$_0^{\mathrm{}}e^{px^2}J_q(ax)J_q(bx)x๐x=\frac{e^{\frac{a^2+b^2}{4p}}}{2p}I_q\left(\frac{ab}{2p}\right)=\frac{e^{\frac{a^2+b^2}{4p}}}{2p}i^qJ_q\left(\frac{iab}{2p}\right)$$
with $`p=i+0^+`$.
We now study the system with the Hamiltonian
$$\begin{array}{cc}\hfill H& =\frac{1}{2}(P_1^2+P_2^2X_1^2X_2^2)=\frac{1}{2}(P^2+\frac{q^2\frac{1}{4}}{X^2}X^2)\hfill \\ & =\frac{1}{2}(S_1U_1+U_1S_1+S_2U_2+U_2S_2)\hfill \\ & =\frac{1}{2}(SP_s+P_sS)\hfill \\ & =\frac{1}{2}(UP_u+P_uU)\hfill \end{array}$$
We will take $`q`$ to be a $`c`$ number and will view the system as having a single degree of freedom. This Hamiltonian has two natural energy eignestates $`|ฯต,in`$ and $`|ฯต,out`$ with wavefunctions and inner products
$$\begin{array}{cc}& s|ฯต,out=\frac{1}{\sqrt{2\pi }}s^{iฯต\frac{1}{2}}\hfill \\ & u|ฯต,in=\frac{1}{\sqrt{2\pi }}u^{iฯต\frac{1}{2}}\hfill \\ & ฯต,out|ฯต^{},in=e^{i\mathrm{\Phi }_A(ฯต)}\delta (ฯตฯต^{})=2^{iฯต}\frac{\mathrm{\Gamma }\left(\frac{1}{2}(1+qiฯต)\right)}{\mathrm{\Gamma }\left(\frac{1}{2}(1+q+iฯต)\right)}\delta (ฯตฯต^{})\hfill \end{array}$$
where in the last inner product we used the integral $`_0^{\mathrm{}}๐yy^{iฯต}J_q(y)=2^{iฯต}\frac{\mathrm{\Gamma }\left(\frac{1}{2}(1+qiฯต)\right)}{\mathrm{\Gamma }\left(\frac{1}{2}(1+q+iฯต)\right)}`$ with nonnegative $`q`$. Note, as a check that these inner products are independent of the sign of $`q`$.
A.2. Second quantized problem
0B
The 0B matrix model is a system of fermions whose first quantized description is the first problem discussed above. Its Lagrangian is
$$=_{\mathrm{}}^{\mathrm{}}๐x\mathrm{\Psi }^{}(x,t)\left(i_t+\frac{1}{2}_x^2+\frac{1}{2}x^2+\mu \right)\mathrm{\Psi }(x,t)$$
In order to express it in terms of the $`s`$ and $`u`$ variables we define the fermionic fields
$$\begin{array}{cc}& \mathrm{\Psi }_s(s,t)=๐xs|x\mathrm{\Psi }(x,t)=\frac{2^{\frac{1}{4}}e^{i\frac{\pi }{8}}}{\sqrt{2\pi }}_{\mathrm{}}^{\mathrm{}}๐x\mathrm{exp}\left(i(\frac{x^2}{2}\sqrt{2}sx+\frac{s^2}{2})\right)\mathrm{\Psi }(x,t)\hfill \\ & \mathrm{\Psi }_u(u,t)=๐xu|x\mathrm{\Psi }(x,t)=\frac{2^{\frac{1}{4}}e^{i\frac{\pi }{8}}}{\sqrt{2\pi }}_{\mathrm{}}^{\mathrm{}}๐x\mathrm{exp}\left(i(\frac{x^2}{2}\sqrt{2}ux\frac{u^2}{2})\right)\mathrm{\Psi }(x,t)\hfill \end{array}$$
which are related through a Fourier transform
$$\mathrm{\Psi }_s(s,t)=๐us|u\mathrm{\Psi }_u(u,t)=\frac{1}{\sqrt{2\pi }}_{\mathrm{}}^{\mathrm{}}๐u\mathrm{exp}\left(isu\right)\mathrm{\Psi }_u(u,t)$$
and find
$$\begin{array}{cc}\hfill & =_{\mathrm{}}^{\mathrm{}}๐s\mathrm{\Psi }_s^{}(s,t)\left(i_t+\frac{i}{2}(s_s+_ss)+\mu \right)\mathrm{\Psi }_s(s,t)\hfill \\ & =_{\mathrm{}}^{\mathrm{}}๐u\mathrm{\Psi }_u^{}(u,t)\left(i_t\frac{i}{2}(u_u+_uu)+\mu \right)\mathrm{\Psi }_u(u,t)\hfill \end{array}$$
We can further simplify the analysis by the change of variables
$$\begin{array}{cc}& \mathrm{\Psi }_1^{(in)}(r,t)=e^{\frac{r}{2}}\mathrm{\Psi }_u(u=e^r,t)\hfill \\ & \mathrm{\Psi }_2^{(in)}(r,t)=e^{\frac{r}{2}}\mathrm{\Psi }_u(u=e^r,t)\hfill \\ & \mathrm{\Psi }_1^{(out)}(r,t)=e^{\frac{r}{2}}\mathrm{\Psi }_s(s=e^r,t)\hfill \\ & \mathrm{\Psi }_2^{(out)}(r,t)=e^{\frac{r}{2}}\mathrm{\Psi }_s(s=e^r,t)\hfill \end{array}$$
which makes the Lagrangians (A.1) look relativistic
$$\begin{array}{cc}\hfill & =_{\mathrm{}}^{\mathrm{}}๐r\underset{i=1,2}{}\mathrm{\Psi }_i^{(in)}(r,t)\left(i_ti_r+\mu \right)\mathrm{\Psi }_i^{(in)}(r,t)\hfill \\ & =_{\mathrm{}}^{\mathrm{}}๐r\underset{i=1,2}{}\mathrm{\Psi }_i^{(out)}(r,t)\left(i_t+i_r+\mu \right)\mathrm{\Psi }_i^{(out)}(r,t)\hfill \end{array}$$
The parameter $`\mu `$ is like the time component of a vector field coupled to the fermion number current whose incoming and outgoing components are
$$\begin{array}{cc}& J^{(in)}=\underset{i}{}\mathrm{\Psi }_i^{(in)}\mathrm{\Psi }_i^{(in)}\hfill \\ & J^{(out)}=\underset{i}{}\mathrm{\Psi }_i^{(out)}\mathrm{\Psi }_i^{(out)}\hfill \end{array}$$
We can remove it by a time dependent gauge transformation, but we prefer not to do so. In this form it is clear that we have four incoming Majorana Weyl fermions $`\mathrm{\Psi }^{(in)}`$ and four outgoing Majorana Weyl fermions $`\mathrm{\Psi }^{(out)}`$ of the opposite chirality. The incoming and the outgoing fermions are related through our map (A.1) which becomes
$$\begin{array}{cc}& \mathrm{\Psi }_1^{(out)}(r,t)=\frac{1}{\sqrt{2\pi }}_{\mathrm{}}^{\mathrm{}}๐r^{}e^{\frac{1}{2}(r+r^{})}\left(\mathrm{exp}(ie^{r+r^{}})\mathrm{\Psi }_1^{(in)}(r^{},t)+\mathrm{exp}(ie^{r+r^{}})\mathrm{\Psi }_2^{(in)}(r^{},t)\right)\hfill \\ & \mathrm{\Psi }_2^{(out)}(r,t)=\frac{1}{\sqrt{2\pi }}_{\mathrm{}}^{\mathrm{}}๐r^{}e^{\frac{1}{2}(r+r^{})}\left(\mathrm{exp}(ie^{r+r^{}})\mathrm{\Psi }_1^{(in)}(r^{},t)+\mathrm{exp}(ie^{r+r^{}})\mathrm{\Psi }_2^{(in)}(r^{},t)\right)\hfill \end{array}$$
What are the symmetries of our problem? In the form (A.1) the Lagrangian has an incoming $`SO(4)SU(2)\times SU(2)^{}`$ symmetry which rotates the incoming fermions and similarly an outgoing $`SO(4)`$ symmetry which rotates the outgoing fermions. The coupling to $`\mu `$ breaks each of these symmetries to $`SU(2)\times U(1)`$. The map between the incoming and outgoing fields (A.1) breaks most of these symmetries. But, let us first ignore the map and start, without loss of generality, by considering the incoming symmetry. The current of the $`U(1)`$ factor has already been mentioned in (A.1). The incoming $`SU(2)`$ currents are $`J^{(in)+}=\mathrm{\Psi }_1^{(in)}\mathrm{\Psi }_2^{(in)}`$, $`J^{(in)}=\mathrm{\Psi }_2^{(in)}\mathrm{\Psi }_1^{(in)}`$ and $`J^{(in)0}=\mathrm{\Psi }_1^{(in)}\mathrm{\Psi }_1^{(in)}\mathrm{\Psi }_2^{(in)}\mathrm{\Psi }_2^{(in)}`$. We note that the currents $`J^{(in)}`$ of (A.1) and $`J^{(in)0}`$ are local in our original โspaceโ coordinate $`s`$, while the currents $`J^{(in)\pm }`$ are nonlocal. The latter involve creating a fermions at $`s`$ and annihilating a fermion at $`s`$, or the other way around. The same distinction between these currents applies in the $`x`$ coordinate.
These four currents have an obvious string theory interpretation. After bosonization the fermion number current $`J^{(in)}`$ creates an incoming NS-NS tachyon $`T^{(in)}`$; roughly<sup>13</sup> We use the word โroughlyโ because of the non-local transform between $`\varphi `$ and $`r`$. $`J^{(in)}(_t_\varphi )T^{(in)}`$. The current $`J^{(in)0}`$ creates the incoming R-R scalar $`C^{(in)}`$; roughly $`J^{(in)0}(_t_\varphi )C^{(in)}`$. These two excitations, which are local in $`x`$, correspond to the perturbative string spectrum. The other two currents $`J^{(in)\pm }`$ create nonperturbative string states. These are solitons โ coherent states of an infinite number of $`C^{(in)}`$ quanta; roughly $`J^{(in)\pm }e^{\pm i\sqrt{2}C^{(in)}}`$. Such solitons were studied in . They create a fermion at one sign of $`x`$ and annihilate it at the other. Note that this is consistent with the $`C`$ field being at the $`SU(2)`$ radius. We now interpret this $`SU(2)`$ symmetry as rotating the two fermion flavors in (A.1).
This discussion of the incoming symmetries is trivially repeated for the outgoing symmetries. In terms of the field $`C`$ the $`SU(2)`$ symmetries of the past and the future are simply those of the left moving and the right moving fields at the selfdual radius.
The map from the past to the future (A.1) shows that only one of the two $`U(1)`$ symmetries is conserved โ the total incoming fermion number equals the total outgoing fermion number. The two $`SU(2)`$ symmetries are more interesting. Both of them are broken, but for large $`|\mu |`$ a certain $`U(1)SU(2)^{(in)}\times SU(2)^{(out)}`$ is approximately conserved. It is broken only by nonperturbative effects of order $`e^{c|\mu |}`$ for some constant $`c`$. The physical interpretation of this fact is simple. We start with a vacuum with $`N\mathrm{}`$ fermions. Let us prepare an initial state with $`\frac{1}{2}N+n^{(in)}`$ incoming fermions from negative $`x`$ and $`\frac{1}{2}Nn^{(in)}`$ incoming fermions from positive $`x`$, and let us examine a final state with $`\frac{1}{2}N+n^{(out)}`$ outgoing fermions to positive $`x`$ and $`\frac{1}{2}Nn^{(out)}`$ outgoing fermions to negative $`x`$. For $`\mu \mathrm{}`$, where the fermions are far below the barrier, there is almost no communication between the left and right sides of the potential, and we must have $`n^{(in)}n^{(out)}`$. Conversely, for $`\mu +\mathrm{}`$ we have $`n^{(in)}n^{(out)}`$. But for finite $`\mu `$ we can have arbitrary $`n^{(in)}`$ and $`n^{(out)}`$. Such scattering processes are created with the insertion of $`n^{(in)}`$ insertions of $`e^{i\sqrt{2}C^{(in)}}`$ in the past (for negative $`n^{(in)}`$ we take $`e^{i\sqrt{2}C^{(in)}}`$), and $`n^{(out)}`$ insertions of $`e^{i\sqrt{2}C^{(out)}}`$ in the future. The condition $`n^{(in)}n^{(out)}`$ for $`\mu +\mathrm{}`$ states that the winding of $`C`$ is approximately conserved while the momentum of $`C`$ is not conserved. For $`\mu \mathrm{}`$ we have the reverse situation.
Let us discuss the discrete symmetries of our problem. First, a $`๐_2`$ subgroup of the the $`SU(2)`$ we mentioned above is not broken by the map (A.1). Combining it with a $`U(1)`$ transformation we can identify it as the parity transformation $`xx`$, $`ss`$ and $`uu`$. In terms of the in and out fermions its action is $`\mathrm{\Psi }_1^{(in/out)}(r,t)\mathrm{\Psi }_2^{(in/out)}(r,t)`$. We identify this transformation with the spacetime charge conjugation which is generated by the worldsheet transformation $`(1)^{F_L}`$ ($`F_L`$ is the leftmoving spacetime fermion number). As a check, note that the currents $`J^{(in/out)}`$ are even and the currents $`J^{(in/out)\pm }`$ are odd under this transformation. Note that spacetime charge conjugation is parity in the matrix model.
Let us consider now the S-duality symmetry in spacetime, which is generated in the worldsheet description by $`(1)^{f_L}`$ ($`f_L`$ is the leftmoving spacetime fermion number). It acts on the parameter $`\mu `$ as $`\mu \mu `$, and therefore it is a symmetry only for $`\mu =0`$. Its action on the fields is
$$\begin{array}{cc}& \mathrm{\Psi }(x,t)\frac{1}{\sqrt{2\pi }}_{\mathrm{}}^{\mathrm{}}๐x^{}e^{ix^{}x}\mathrm{\Psi }^{}(x^{},t)\hfill \\ & \mathrm{\Psi }_s(s,t)\mathrm{\Psi }_s^{}(s,t)\hfill \\ & \mathrm{\Psi }_u(u,t)\mathrm{\Psi }_u^{}(u,t)\hfill \\ & \mathrm{\Psi }_1^{(in)}\mathrm{\Psi }_2^{(in)}\hfill \\ & \mathrm{\Psi }_i^{(out)}\mathrm{\Psi }_i^{(out)}\hfill \end{array}$$
It is easy to check that it is a symmetry of the Lagrangian (A.1)(A.1)(A.1) with $`\mu =0`$, and of the transforms (A.1)(A.1)(A.1). Clearly, if we combine this operation with the parity transformation $`(1)^{F_L}`$ the difference between the transformations of $`\mathrm{\Psi }_s`$ and $`\mathrm{\Psi }_u`$ and the difference between $`\mathrm{\Psi }^{(in)}`$ and $`\mathrm{\Psi }^{(out)}`$ are reversed. It is interesting that while in $`x`$ space the transformation involves a duality transformation of $`(x,p)`$, which is implemented by a Fourier transform, in $`s`$, $`u`$ and $`r`$ space no such duality transformation is needed, and the transformation rules are local and simple. We conclude that in the $`s`$, $`u`$ and $`r`$ variables this $`๐_2`$ symmetry acts as charge conjugation on the outgoing fields and as $`CP`$ on the ingoing fields.
We point out that this $`๐_2`$ symmetry is a subgroup of the original $`SO(4)`$ we mentioned above.
0A
The 0A matrix model is a system of fermions whose first quantized description is the second problem discussed above. Its Lagrangian is
$$=_0^{\mathrm{}}๐x\mathrm{\Psi }^{}(x,t)\left(i_t+\frac{1}{2}(_x^2+x^2\frac{q^2\frac{1}{4}}{x^2})+\mu \right)\mathrm{\Psi }(x,t)$$
In order to express it in terms of the $`s`$ and $`u`$ variables we define new fermionic fields which are related by integral transforms
$$\begin{array}{cc}& \mathrm{\Psi }_s(s,t)=๐xs,q|x,q\mathrm{\Psi }(x,t)=e^{i\pi (\frac{1}{4}\frac{3q}{4})}_0^{\mathrm{}}๐x\sqrt{2xs}e^{\frac{i}{2}(x^2+s^2)}J_q(\sqrt{2}sx)\mathrm{\Psi }(x,t)\hfill \\ & \mathrm{\Psi }_u(u,t)=๐xu,q|x,q\mathrm{\Psi }(x,t)=e^{i\pi (\frac{1}{4}\frac{3q}{4})}_0^{\mathrm{}}๐x\sqrt{2xu}e^{\frac{i}{2}(x^2+u^2)}J_q(\sqrt{2}ux)\mathrm{\Psi }(x,t)\hfill \\ & \mathrm{\Psi }_s(s,t)=๐us,q|u,q\mathrm{\Psi }_u(u,t)=_0^{\mathrm{}}๐u\sqrt{su}J_q(su)\mathrm{\Psi }_u(u,t)\hfill \end{array}$$
As in (A.1) it is convenient to express them in terms of incoming and outgoing fermions
$$\begin{array}{cc}& \mathrm{\Psi }^{(in)}(r,t)=e^{\frac{r}{2}}\mathrm{\Psi }_u(u=e^r)\hfill \\ & \mathrm{\Psi }^{(out)}(r,t)=e^{\frac{r}{2}}\mathrm{\Psi }_s(s=e^r)\hfill \\ & \mathrm{\Psi }^{(out)}(r,t)=_{\mathrm{}}^{\mathrm{}}๐r^{}e^{r+r^{}}J_q(e^{r+r^{}})\mathrm{\Psi }^{(in)}(r^{},t)\hfill \end{array}$$
It is easy to express the Lagrangian (A.1) in these variables
$$\begin{array}{cc}\hfill & =_0^{\mathrm{}}๐s\mathrm{\Psi }_s^{}(s,t)\left(i_t+\frac{i}{2}(s_s+_ss)+\mu \right)\mathrm{\Psi }_s(s,t)\hfill \\ & =_0^{\mathrm{}}๐u\mathrm{\Psi }_u^{}(u,t)\left(i_t\frac{i}{2}(u_u+_uu)+\mu \right)\mathrm{\Psi }_u(u,t)\hfill \\ & =_0^{\mathrm{}}๐r\mathrm{\Psi }^{(in)}(r,t)\left(i_ti_r+\mu \right)\mathrm{\Psi }^{(in)}(r,t)\hfill \\ & =_0^{\mathrm{}}๐r\mathrm{\Psi }^{(out)}(r,t)\left(i_t+i_r+\mu \right)\mathrm{\Psi }^{(out)}(r,t)\hfill \end{array}$$
As in the 0B theory we have a $`U(1)`$ symmetry which rotates $`\mathrm{\Psi }`$ by a phase and corresponds to fermion number conservation. The spacetime charge conjugation symmetry, which is generated by $`(1)^{F_L}`$ on the worldsheet, acts as $`qq`$. Up to an overall phase this acts leaving $`\mathrm{\Psi }^{(in)}`$ invariant and changing $`\mathrm{\Psi }^{(out)}(1)^q\mathrm{\Psi }^{(out)}`$, which is somewhat trivial. The S-duality transformation, $`(1)^{f_L}`$, acts on the parameter $`\mu `$ as $`\mu \mu `$, and is a symmetry only for $`\mu =0`$. Its action on the fields is
$$\begin{array}{cc}& \mathrm{\Psi }(x,t)_0^{\mathrm{}}๐x^{}\sqrt{xx^{}}J_q(xx^{})\mathrm{\Psi }^{}(x^{},t)\hfill \\ & \mathrm{\Psi }_s(s,t)\mathrm{\Psi }_s^{}(s,t)\hfill \\ & \mathrm{\Psi }_u(u,t)\mathrm{\Psi }_u^{}(u,t)\hfill \\ & \mathrm{\Psi }^{(in)}(r,t)\mathrm{\Psi }^{(in)}(r,t)\hfill \\ & \mathrm{\Psi }^{(out)}(r,t)\mathrm{\Psi }^{(out)}(r,t)\hfill \end{array}$$
It is easy to check that it is a symmetry of the Lagrangian (A.1)(A.1) with $`\mu =0`$, and of the transforms (A.1). It is interesting that while in $`x`$ space the transformation involves an integral transform, in $`s`$ and $`u`$ space the transformation is local and it is very simple. We conclude that the S-duality symmetry acts as charge conjugation of the matrix model fermions in the $`s`$, $`u`$ and $`r`$ variables.
A.3. Computation of the free energies
We are interested in computing the partition function of the thermal system. Using the density of states $`\rho (ฯต)=\frac{\varphi ^{}(ฯต)}{2\pi }`$ we can write the standard expression is
$$\begin{array}{cc}\hfill \mathrm{log}๐ต=& \underset{\mathrm{\Lambda }\mathrm{}}{lim}\frac{1}{2\pi }_\mathrm{\Lambda }^{\mathrm{}}๐ฯต\varphi ^{}(ฯต)\mathrm{log}(1+e^{2\pi R(ฯต\mu )})\hfill \\ \hfill =& \underset{\mathrm{\Lambda }\mathrm{}}{lim}R\left(\varphi (\mathrm{\Lambda })(\mathrm{\Lambda }+\mu )+_\mathrm{\Lambda }^{\mathrm{}}๐ฯต\frac{\varphi (ฯต)}{1+e^{2\pi R(ฯต\mu )}}\right)\hfill \end{array}$$
where $`\mathrm{\Lambda }`$ is a cutoff on the bottom of the Fermi sea and we neglected terms which are exponentially small at large $`\mathrm{\Lambda }`$. In order to simplify the analysis and not worrying about the $`\mathrm{\Lambda }`$ dependence, we will study the second derivative of (A.1)
$$_\mu ^2\mathrm{log}๐ต=R_{\mathrm{}}^{\mathrm{}}๐ฯต_\mu ^2\frac{\varphi (ฯต)}{1+e^{2\pi R(ฯต\mu )}}$$
which is a convergent integral.
Let us now consider the 0B theory. The determinant of the single particle S-matrix of the 0B theory is $`ie^{i\mathrm{\Phi }_B(ฯต)}`$. It can be expressed as
$$\mathrm{\Phi }_B(ฯต)=i\mathrm{log}\left(\frac{\mathrm{\Gamma }(\frac{1}{2}iฯต)}{\mathrm{\Gamma }(\frac{1}{2}+iฯต)}\right)=_0^{\mathrm{}}\frac{dt}{t}\left(\frac{\mathrm{sin}(ฯตt)}{\mathrm{sinh}\frac{t}{2}}2ฯตe^t\right)$$
where we have used (4.1). Then, (A.1) becomes
$$\begin{array}{cc}\hfill _\mu ^2\mathrm{log}๐ต_B=& R_0^{\mathrm{}}\frac{dt}{t}_{\mathrm{}}^{\mathrm{}}๐ฯต_\mu ^2\frac{1}{1+e^{2\pi R(ฯต\mu )}}\left(\frac{\mathrm{sin}(ฯตt)}{\mathrm{sinh}\frac{t}{2}}2ฯตe^t\right)\hfill \\ \hfill =& _0^{\mathrm{}}\frac{dt}{t}_\mu ^2\left(\frac{\mathrm{cos}(\mu t)}{2\mathrm{sinh}\frac{t}{2}\mathrm{sinh}\frac{t}{2R}}+R\mu ^2e^t\right)\hfill \end{array}$$
We now want to integrate this equation twice with respect to $`\mu `$. Invariance under $`\mu \mu `$ forbids a term linear in $`\mu `$ and the constant term is fixed arbitrarily such that
$$\begin{array}{cc}\hfill \mathrm{log}๐ต_B=& _0^{\mathrm{}}\frac{dt}{t}\left(\frac{\mathrm{cos}(\mu t)}{2\mathrm{sinh}\frac{t}{2}\mathrm{sinh}\frac{t}{2R}}\frac{2R}{t^2}+\left[\frac{1}{12}(R+\frac{1}{R})+R\mu ^2\right]e^t\right)\hfill \\ \hfill =& \mathrm{\Omega }(y=2i\mu ,R)+\mathrm{\Omega }(\overline{y}=2i\mu ,R)\hfill \end{array}$$
The single particle S-matrix in the 0A theory with nonzero $`q`$ is $`e^{i\mathrm{\Phi }_A(ฯต)}`$. It can be expressed as
$$\mathrm{\Phi }_A(ฯต)=i\mathrm{log}\left(\frac{\mathrm{\Gamma }(\frac{1}{2}+\frac{1}{2}(qiฯต))}{\mathrm{\Gamma }(\frac{1}{2}+\frac{1}{2}(q+iฯต))}\right)=_0^{\mathrm{}}\frac{dt}{t}\left(\frac{e^{qt/2}\mathrm{sin}(\frac{ฯตt}{2})}{\mathrm{sinh}\frac{t}{2}}ฯตe^t\right)$$
where we have used (4.1), and we dropped a constant term in the phase, as well as a term that is linear in $`ฯต`$. The term linear in $`ฯต`$ could be removed by doing a rescaling of the variables $`u`$ and $`s`$ that appeared in the 0A discussion<sup>14</sup> This term that is linear in $`ฯต`$ would have lead to an extra term proportional to $`\mu ^2`$ in the free energy. We choose to remove this analytic term by hand. Then, (A.1) becomes
$$\begin{array}{cc}\hfill _\mu ^2\mathrm{log}๐ต_A=& R_0^{\mathrm{}}\frac{dt}{t}_{\mathrm{}}^{\mathrm{}}๐ฯต_\mu ^2\frac{1}{1+e^{2\pi R(ฯต\mu )}}\left(\frac{e^{qt/2}\mathrm{sin}(\frac{ฯตt}{2})}{\mathrm{sinh}\frac{t}{2}}ฯตe^t\right)\hfill \\ \hfill =& _0^{\mathrm{}}\frac{dt}{t}_\mu ^2\left(\frac{e^{qt/2}\mathrm{cos}(\frac{\mu t}{2})}{2\mathrm{sinh}\frac{t}{2}\mathrm{sinh}\frac{t}{4R}}+\frac{R\mu ^2}{2}e^t\right)\hfill \\ \hfill =& _\mu ^2\left[\mathrm{\Omega }(y=q+i\mu ,2R)+\mathrm{\Omega }(\overline{y}=qi\mu ,2R)\right]\hfill \end{array}$$
Using the definition of the function $`\mathrm{\Omega }(y,r)`$ (3.1) one can derive
$$\mathrm{\Omega }(yr,\frac{1}{r})=\mathrm{\Omega }(y,r)[\frac{1}{24}(r+\frac{1}{r})\frac{ry^2}{8}]\mathrm{log}r$$
Using this relation we can check that the 0B answer (A.1) is the same as the 0A answer (A.1) for $`q=0`$, up to a term involving $`\mathrm{log}r`$. This term arises because in 0A and 0B it is natural to choose the cutoffs to be $`R`$ independent. On the other hand T-duality relates them by $`\mathrm{\Lambda }_B=\mathrm{\Lambda }_A\frac{R_A}{\sqrt{2\alpha ^{}}}`$ (see (1.1)). Once we take this into account, the terms that are logarithmic in the cutoff give a contribution cancelling the last term in (A.1).
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# Numerical Ricci-flat metrics on K3
## 1 Introduction
In 1977, S.-T. Yau proved E. Calabiโs conjecture that a compact Kรคhler manifold with vanishing first Chern class admits a Ricci-flat metric in each Kรคhler class. Interest in Calabi-Yau manifolds was subsequently generated among physicists by the discovery that they can serve as supersymmetry-preserving compactification manifolds in string theory . In the three decades since the proof of Yauโs theorem, many examples of Calabi-Yau manifolds have been constructed and studied, principally using methods of algebraic geometry, and much has been learned about their mathematical properties and physical applications. Nonetheless, a major gap in our knowledge about Calabi-Yau manifolds has persisted during this time, namely the Ricci-flat metrics themselves. Yauโs proof is not constructive, and no example of a smooth Ricci-flat metric is known explicitly for any Calabi-Yau manifold. Indeed, perhaps we should not expect there to exist such metrics in closed form .
The question thus arises as to whether it is possible to solve the Einstein equation numerically on a Calabi-Yau manifold. The purpose of this paper it to show that it is. We describe a method for doing so, and display the metrics obtained by applying that method to the smallest-dimensional Calabi-Yau, the K3 surface.
The algorithms we developed rely in an essential way on the underlying complex and Kรคhler geometry of Calabi-Yau manifolds. In fact, one of the main points of this paper is to show that those properties are as powerful for numerical work as they have already proven to be for analytical calculations and proving theorems. To explain this, let us consider the challenges faced by someone attempting to numerically solve the Euclidean Einstein equation on a general real four-manifold. In some sense this is a problem in numerical relativity, and one can get a sense of its scale by considering that four-dimensional problems in numerical relativity, such as black hole collisions, can be solved (if at all) only with the investment of extremely large computing resources, typically supercomputers. (Note that Calabi-Yau manifolds admit no continuous isometries that could reduce the effective dimensionality of the problem.) At a more fundamental level, whereas conventional numerical relativity deals with solving the Einstein equation as an initial-value problem on a Lorentzian spacetime, here we wish to solve it on a Euclidean manifold, a problem for which general algorithms are lacking. The challenges for creating an algorithm include the usual issues of gauge fixing and coordinate singularities, as well as finding a way to fix any moduli the solutions might have. An even more difficult challenge would be to avoid (or deal with) the curvature singularities that generically form under relaxation schemes such as Ricci flow.
Let us now see how the framework of complex and Kรคhler differential geometry allows one to naturally solve, or greatly ameliorate, each of these problems in turn. Firstly, the Kรคhler formulation of geometry can be employed to vastly reduce the scale of the problem compared to the language of real differential geometry. The metric can be encoded in a single scalar function, the Kรคhler potential (as reviewed in subsection 2.1). In terms of the Kรคhler potential the Einstein equation takes the relatively simple form of a Monge-Ampรจre equation (as reviewed in subsection 2.2). So we have a large simplication of the many degrees of freedom required to locally parameterize a general metric, and of the complicated set of differential equations in the usual form of the Einstein equation.
Secondly, complex coordinates offer a naturally adapted gauge choice, since with respect to any Kรคhler metric they satisfy the harmonic gauge condition, a well-known gauge choice for numerical relativity. Furthermore the gauge fixing is nearly complete: there are only a finite number of (continuous) residual gauge transformations, since compact complex manifolds admit only a finite number of holomorphic vector fields. In fact, for Calabi-Yaus this number is zero, so the gauge fixing is complete. Since there are no gauge transformations, no coordinate singularities can appear.
Thirdly, the moduli of Calabi-Yau manifolds, which are divided into complex structure and Kรคhler moduli, may be fixed at any desired valuesโbefore solving the Einstein equationโin the following way. With the manifold defined topologically by an atlas of patches, the complex structure is fixed by fixing the holomorphic coordinate transition functions on the patch overlaps, while the Kรคhler moduli are fixed by fixing the Kรคhler transformations on them (which then serve as the boundary conditions for the Kรคhler potential). (Details of the above are given in subsection 2.1.)
Finally, there is a kind of stability that appears to be inherent in Kรคhler geometry, miraculously eliminating the problem of spontaneous formation of curvature singularities. Ricci flow, for example, contrary to its behaviour on real manifolds, is extremely robust on Kรคhler manifolds, as shown by Caoโs long-time existence theorems . On Calabi-Yau manifolds in particular, starting from *any* Kรคhler metric it converges to the Ricci-flat metric in the same class. In principle therefore it provides a general algorithm for solving the Einstein equation on Calabi-Yaus. However, since it is a rather inefficient method from a computational viewpoint (rather like solving the Laplace equation by simulating diffusion), we developed and used instead a Gauss-Seidel-type relaxation algorithm for the Monge-Ampรจre equation. While we lack a convergence theorem for our algorithm, we found that in practice it was just as robust as Ricci flow, presumably for the same underlying reasons. (In fact, in one sense our algorithm was even more robust than Ricci flow, since remarkably it converged even when the initial Kรคhler potential didnโt define a \[positive-definite\] metric.) Subsection 2.3 contains a discussion of Ricci flow and a description of our algorithm.
As a proof of principle for our methods, we applied them to a class of K3โs known as Kummer surfaces, which are blow-ups of the orbifold $`T^4/๐_2`$. For simplicity, we considered the most symmetrical Kummer surfaces, namely those for which the torus is cubical and all 16 fixed points are blown up identically. This symmetry leaves only one modulus (not counting the trivial volume modulus), namely the ratio of the size of the blow-ups to the size of the $`T^4`$. In fact there is a limit to how large this ratio may be, since at some point certain holomorphic curves shrink to zero size, signalling the appearance of new orbifold singularitiesโon this wall of the Kรคhler cone the manifold is an orbifold of another, smooth K3. (Details of the construction are given in subsection 3.1.) Using our algorithm, we computed the Ricci-flat metric at various points over the full range of this modulus. Each point required a few days on a garden-variety desktop computer for the highest resolutions. Subsection 3.2 is devoted to an exploration of the resulting geometries as a function of the modulus. Various curvature invariants are plotted, as well as a low-lying eigenvalue of the scalar Laplacian.
The success we had with these highly symmetrical K3 surfaces leads us to consider possible generalizations. Using our methods, could we solve the Einstein equation on generic K3โs, which lack such discrete symmetries, or on Calabi-Yau three-folds? What about on other Kรคhler manifolds such as del Pezzo surfaces, or on Calabi-Yaus with matter such as fluxes and branes? The estimates we make in subsection 4.1 suggest that the answer to these questions is yes, but that due to memory limitations in the three-fold case we would need to be helped by a high degree of discrete symmetry.
In subsection 4.2 we return to the problem of solving the Einstein equation in the real Euclidean context, and explore what this work has taught us that might generalize to non-Kรคhler geometries. We conclude in subsection 4.3 by mentioning some possible mathematical and physical applications of these numerical metrics.
The C code for our simulations, as well as animated versions of the plots shown in this paper, are available at the website http://schwinger.harvard.edu/~wiseman/K3/.
## 2 General method
The purpose of this section is to discuss in a general way the problem of solving the Einstein equation numerically on a Calabi-Yau manifold at a given point in its moduli space. We first briefly review the essentials of Kรคhler geometry,<sup>1</sup><sup>1</sup>1A detailed review of Kรคhler geometry and Calabi-Yau manifolds may be found in the excellent set of lecture notes . describing how the geometry is encoded in the Kรคhler potential and how the moduli are fixed. We then explain how, in Kรคhler geometry, the Einstein equation reduces to a Monge-Ampรจre equation. Finally, we discuss in general terms the numerical algorithms we applied to solving that equation.
### 2.1 Kรคhler geometry
We work on a manifold with a fixed complex structure, that is, on each coordinate patch $`U_\alpha `$ we have a set of complex coordinates $`\{z_\alpha ^i,\overline{z}_\alpha ^i\}`$ such that on each overlap $`U_\alpha U_\beta `$ the transition functions $`z_\beta ^i(z_\alpha )`$ are holomorphic ($`\alpha ,\beta `$ index the patch and $`i,j`$ the complex coordinate). Fixing the complex structure restricts the allowed gauge transformations to holomorphic diffeomorphisms. Continuous holomorphic diffeomorphisms are generated by holomorphic vector fields. In contrast to the infinite number of vector fields generating general diffeomorphisms, it can be shown that a compact manifold admits only a finite number of holomorphic vector fields. Even better, Calabi-Yau manifolds have none at all, for the following reason. Every Calabi-Yau is equipped with a nowhere-vanishing holomorphic $`(n,0)`$ form $`\mathrm{\Omega }`$ (which depends on the complex structure but not on the metric). If $`v^i`$ is non-vanishing with holomorphic components, then the same is true of the $`(n1,0)`$ form $`v^{i_1}\mathrm{\Omega }_{i_1i_2\mathrm{}i_n}dz^{i_2}\mathrm{}dz^{i_n}`$. Since the Hodge number $`h^{n1,0}`$ vanishes on a Calabi-Yau, such a form must be $`\overline{}`$ exact, which is impossible for an $`(n1,0)`$ form. Hence, on a Calabi-Yau, fixing the transition functions for the complex coordinates amounts to a full gauge fixing. As mentioned in the introduction, this is a great advantage for numerical work.
A metric is called Kรคhler with respect to a given complex structure if it is Hermitian, $`g_{ij}=g_{\overline{ฤฑ}\overline{ศท}}=0`$, and if the associated Kรคhler form $`J=ig_{i\overline{ศท}}dz^id\overline{z}^j`$ is closed, $`dJ=0`$. $`J`$ is obviously a positive $`(1,1)`$ form, since the metric is everywhere positive definite. The cohomology class of $`J`$ is called the Kรคhler class. The set of potential Kรคhler classes, i.e. $`(1,1)`$ cohomology classes containing at least one positive form, is called the Kรคhler cone, or Kรคhler moduli space.
The fact that the Kรคhler form is closed implies that it, and the metric, are locally expressible as the matrix of second derivatives of a scalar,
$$g_{i\overline{ศท}}|_{U_\alpha }=_i_{\overline{ศท}}K_\alpha ,J|_{U_\alpha }=i\overline{}K_\alpha ,$$
(1)
where $`K_\alpha `$ is a real function, the Kรคhler potential, defined on the patch $`U_\alpha `$. Given $`g_{i\overline{ศท}}`$, the Kรคhler potential is not unique but can be changed by adding the real part of any holomorphic function, a so-called Kรคhler transformation. From a numerical point of view, one advantage of using the Kรคhler potential to encode the geometry is obvious: it reduces the number of functions to store from $`n(2n+1)`$ for the full metric (or $`n^2`$ if the hermiticity condition is imposed) to a single one.
The volume of the manifold is given in terms of the Kรคhler form by
$$V=\frac{1}{n!}J^n.$$
(2)
Therefore on a compact manifold $`J`$ cannot be exact, or else the manifold would have zero volume. So the Kรคhler potentials obtained by integrating (1) must disagree on the patch overlaps by Kรคhler transformations:
$$K_\alpha K_\beta =u_{\alpha \beta }.$$
(3)
These Kรคhler transformations serve as boundary conditions for the Kรคhler potential (the only boundary conditions if the manifold is compact). In doing so, they perform two important tasks. First, they fix the Kรคhler class. To see this, note that two Kรคhler potentials $`\{K_\alpha \}`$ and $`\{K_\alpha ^{}\}`$ which have the same Kรคhler transformations, $`K_\alpha K_\beta =K_\alpha ^{}K_\beta ^{}`$, differ by a globally defined real function $`\varphi `$. Therefore the difference between the corresponding Kรคhler forms, $`JJ^{}=i\overline{}\varphi `$, is an exact form. Conversely, given a representative of some particular Kรคhler class, a corresponding set of $`u_{\alpha \beta }`$โs may easily be found by solving (1) separately on each patch. (The representative need not be positive. This is useful particularly near the edge of the Kรคhler cone, where it may not be easy to find positive representatives.) The second task performed by the Kรคhler transformations is to (almost) eliminate gauge-equivalent Kรคhler potentials, that is, different Kรคhler potentials that give rise to the same metric: for each Kรคhler form $`J`$ in the class defined by a set $`\{u_{\alpha \beta }\}`$, there is a unique solution to (1) and (3), up to constant shifts of all the $`K_\alpha `$.
When the metric is Kรคhler, only the purely holomorphic and antiholomorphic components $`\mathrm{\Gamma }_{jk}^i`$ and $`\mathrm{\Gamma }_{\overline{ศท}\overline{k}}^{\overline{ฤฑ}}`$ of the Christoffel symbol are non-zero. Together with the hermiticity of the metric, it follows that the coordinates $`z^i`$ and $`\overline{z}^i`$ are harmonic, a well-known gauge fixing condition for numerical relativity. However, Kรคhlerity is a stronger condition than harmonicity; whereas a given metric always admits a harmonic coordinate system, at least locally, a generic metric is not Kรคhler with respect to any coordinates, even locally.<sup>2</sup><sup>2</sup>2Nor are Kรคhler coordinates necessarily unique for a given metric. For example, a Ricci-flat metric on a hyperkรคhler manifold, such as K3, is Kรคhler with respect to a continuous family of different complex structures. We will not make use of this additional structure, but will be content to fix a particular complex structure at the outset.
To summarize: We fix the complex structure by specifying the coordinate transition functions on the patch overlaps. Using the Kรคhler potential to encode the geometry, we fix the Kรคhler class by specifying the Kรคhler transformations on the overlaps. With this set-up in hand, we now turn to the question of what equation to solve for the Kรคhler potential.
### 2.2 The Monge-Ampรจre equation
There exists a simple expression for the Ricci tensor of a Kรคhler metric:
$$R_{kl}=R_{\overline{k}\overline{l}}=0,R_{k\overline{l}}=R_{\overline{l}k}=_k_{\overline{l}}\mathrm{ln}detg_{i\overline{ศท}}$$
(4)
(note that $`detg_{i\overline{ศท}}=\sqrt{|g|}`$). Just as we defined the Kรคhler form $`J`$, we can define the Ricci form $`=iR_{k\overline{l}}dz^kd\overline{z}^l`$. By virtue of (4), $``$ is closed (but not necessarily exact, since $`detg_{i\overline{ศท}}`$ is not in general a globally defined function). One can show that its cohomology class, the first Chern class $`c_1`$, is a topological invariant: two different Kรคhler metrics will give rise to Ricci forms that differ by an exact $`(1,1)`$ form. Calabi-Yau manifolds are defined by the condition that $`c_1`$ vanishes, in other words that the Ricci form is always exact. According to Yauโs theorem , on a Calabi-Yau manifold each Kรคhler class (i.e. $`(1,1)`$ cohomology class containing at least one Kรคhler form) contains a unique Ricci-flat Kรคhler form.
From (4), the Einstein equation for a Kรคhler metric reads
$$_k_{\overline{l}}\mathrm{ln}detg_{i\overline{ศท}}=0.$$
(5)
In terms of the metric this is a second-order PDE, but in terms of the Kรคhler potential it is fourth order. It might seem then that working with the Kรคhler potential is not so advantageous after all. However, as we will now explain, by the magic of complex analysis it can be reduced to a second-order PDE, specifically a Monge-Ampรจre equation (a PDE in which the derivatives appear in the form of a Hessian).
For simplicity, let us assume first that the coordinates have been arranged in such a way that the Jacobians $`det_{ij}(z_\alpha ^i/z_\beta ^j)`$ of the transition functions are 1 on all the overlaps. In that case $`detg_{i\overline{ศท}}`$ is a globally defined function, and on a compact manifold the Einstein equation (5) is equivalent to it being constant:
$$detg_{i\overline{ศท}}=\lambda .$$
(6)
This is a non-linear Monge-Ampรจre equation for the Kรคhler potential. The constant $`\lambda `$ is related to the volume of the manifold, given by (2), which depends only on the Kรคhler class and may therefore be calculated a priori from the Kรคhler transformations on the overlaps (see Appendix B for details).
The coordinate system used for the Kummer surface in Section 3 happens to satisfy the condition of unit coordinate Jacobians assumed in the previous paragraph. It is instructive nonetheless to consider what one would do in the more general situation. Clearly we need to replace the right-hand side of (6) with something that has the correct transformation law on the overlaps, and also that implies the Einstein equation (5). The holomorphic $`(n,0)`$ form $`\mathrm{\Omega }`$ once again comes to the rescue; the generalization we seek is:
$$detg_{i\overline{ศท}}=\lambda |\mathrm{\Omega }_{1\mathrm{}n}|^2$$
(7)
(this can also be written in the coordinate-invariant form $`J^n=(i)^nn!\lambda \mathrm{\Omega }\overline{\mathrm{\Omega }}`$). If $`\mathrm{\Omega }`$ is not known explicitly, then $`|\mathrm{\Omega }_{1\mathrm{}n}|^2`$ can be found as follows in terms of an arbitrary Kรคhler metric $`\stackrel{~}{g}_{i\overline{ศท}}`$ (which need not be in the same Kรคhler class as the desired solution, since $`\mathrm{\Omega }`$ depends only on the complex structure). Since the manifold is Calabi-Yau, its Ricci form is exact, $`\stackrel{~}{R}_{i\overline{ศท}}=_i_{\overline{ศท}}\stackrel{~}{F}`$, where $`\stackrel{~}{F}`$ is a globally defined function which can be calculated explicitly (either analytically or by standard numerical methods). We then have
$$|\mathrm{\Omega }_{1\mathrm{}n}|^2=e^{\stackrel{~}{F}}det\stackrel{~}{g}_{i\overline{ศท}},$$
(8)
and equation (7) may be applied.
To get a feeling for the Monge-Ampรจre equation, it is useful to consider equations (7) and (8) in the case where $`\stackrel{~}{g}_{i\overline{ศท}}`$ *is* in the desired Kรคhler class. Then we can write
$$K_\alpha =\stackrel{~}{K}_\alpha +\varphi ,$$
(9)
where $`\varphi `$ is a globally defined scalar. In terms of $`\varphi `$ the Monge-Ampรจre equation is
$$det(\delta _i^k+\stackrel{~}{g}^{k\overline{ศท}}_i_{\overline{ศท}}\varphi )=\lambda e^{\stackrel{~}{F}}.$$
(10)
The left-hand side is a non-linear operator acting on $`\varphi `$. If $`_i_{\overline{ศท}}\varphi `$ is small (i.e. if $`\stackrel{~}{g}_{i\overline{ศท}}`$ is almost Ricci-flat) then we can linearize, yielding a Poisson equation:
$$\frac{1}{2}\stackrel{~}{}^2\varphi +๐ช\left(_i_{\overline{ศท}}\varphi \right)^2=\lambda e^{\stackrel{~}{F}}1.$$
(11)
In this sense Yauโs theorem can be understood as a generalization of the existence theorem for the Poisson equation.
### 2.3 Methods
The similarity to the Poisson equation observed in the last subsection suggests that methods for solving it might be generalized to the Monge-Ampรจre equation. Among the standard methods for the Poisson equation, some are local and others are non-local. We have restricted ourselves to local schemes, for two reasons. First, the Monge-Ampรจre equation is local and non-linear. Second, for most manifolds (even highly symmetrical ones), the individual patches $`U_\alpha `$ (or, more precisely, their images in $`๐^n`$, i.e. the ranges of the coordinates $`z_\alpha ^i`$), have rather irregular shapes. This motivated the use of a lattice discretization, rather than non-local spectral representations such as Fourier modes or wavelets.
One simple (but inefficient) method for solving the Poisson equation is to simulate diffusion. The analogous geometric relaxation equation is Ricci flow, which is defined by a first order equation in an auxiliary time dimension: $`\dot{g}_{i\overline{ศท}}=R_{i\overline{ศท}}`$. This flow is well studied mathematically , largely because of its application to geometrization. It is also of some interest in physics as the one-loop renormalization-group evolution for the target space geometry of a sigma model. In terms of the Kรคhler potential the flow is governed by
$$\dot{K}_\alpha =\mathrm{ln}\left(\frac{detg_{i\overline{ศท}}}{\lambda |\mathrm{\Omega }_{1\mathrm{}n}|^2}\right).$$
(12)
Note that on the overlaps $`\dot{K}_\alpha \dot{K}_\beta =0`$, so the Kรคhler transformations, and hence the Kรคhler moduli, are conserved. Cao has shown that Ricci flow starting from an arbitrary Kรคhler metric converges to the Ricci-flat metric in its class.
Some work has been done on numerical simulation of Ricci flow on real geometries , but not (as far as we know) in the Kรคhler case. We experimented with simple lattice implementations of such a scheme. However, if one is only interested in the endpoint of the flow, namely the Ricci-flat metric, as opposed to the whole flow, this method is clearly very inefficient, as it requires solving an equation in one higher dimension. Furthermore, stability of the diffusion problem typically requires implicit finite differencing schemes which are rather inconvenient (particularly in several dimensions). Another more subtle drawback which nonetheless is rather serious in practice is that it requires an initial Kรคhler (i.e. positive) form in the desired class, as seen above explicitly from the logarithm in (12). As mentioned in subsection 2.1, it is not always easy to find a positive representative of a given class, especially near the edge of the Kรคhler cone.
The prototype local method for solving the Poisson equation is the Gauss-Seidel method. Here the discretized Poisson equation is solved at each lattice point in turn. After a suitable number of iterations over the whole lattice, the discretized equations will be solved to a given accuracy. This is very robust and simple to implement. Furthermore whilst Gauss-Seidel is slow compared to spectral or multigrid methods in low dimension, scaling as $`N^{1+2/d}`$ rather than $`N\mathrm{log}N`$ in $`d`$ real dimensions (where $`N`$ is the number of lattice points), in our K3 case of 4 real dimensions, and more so for still larger dimensions, the advantage is not so great. Of course multigrid could be implemented relatively simply to improve our Gauss-Seidel method if speed became a crucial issue.
Our analog of the Gauss-Seidel method for the Monge-Ampรจre equation is as follows. On a lattice the metric is determined from the Kรคhler potential by taking discrete derivatives. Thus the value of $`detg_{i\overline{ศท}}`$ at a given site is a function of the values of $`K_\alpha `$ at that site and its neighbors, out to some distance depending on the order of finite differencing used. Our algorithm directs one to go through each site of the lattice in turn, changing $`K_\alpha `$ at that site to the value which solves (7) given its values at the neighboring sites. Although we do not have a convergence theorem for this algorithm, we found that in practice it converged on the full range of Kรคhler parameters studied. This included cases where the initial Kรคhler potential did *not* define a positive $`(1,1)`$ form.
As noted earlier, the constant $`\lambda `$ which relates the coordinate volume to the proper volume can be computed analytically from the Kรคhler class. Now consider solving the equation
$$detg_{i\overline{ศท}}=\stackrel{~}{\lambda }.$$
(13)
We will *only* find a solution to this new Monge-Ampรจre equation when we set $`\stackrel{~}{\lambda }=\lambda `$. However, there is a subtlety, as when we implement this equation numerically, the errors involved in finite differencing will mean that the value of $`\stackrel{~}{\lambda }`$ that solves the finite differenced equation will actually differ slightly from the true continuum value $`\lambda `$. This implies the finite differenced equation should have no solution when $`\stackrel{~}{\lambda }=\lambda `$, which sounds mildly disastrous. However, we can see in detail how this discretization error manifests itself if we consider the Ricci flow equation (12) with $`\lambda `$ replaced by $`\stackrel{~}{\lambda }`$. Then, since the Monge-Ampรฉre equation has no solution for $`\stackrel{~}{\lambda }\lambda `$, the flow will never reach a fixed point. However, the flow does asymptote to one with a simple time dependence, namely,
$$K_{\stackrel{~}{\lambda }}(t)=K_\lambda +tv,v=\mathrm{ln}\frac{\lambda }{\stackrel{~}{\lambda }},$$
(14)
where $`K_\lambda `$ is the solution of the Monge-Ampรจre equation for the true value of $`\lambda `$. The constant $`v`$, as determined above, then gives the asymptotic time dependence of the flow. By analogy with the Ricci flow above, this implies that under Gauss-Seidel iteration, our finite differenced Monge-Ampรจre equation (13) will have an asymptotic solution that drifts in iteration time by a constant mode. This simply corresponds to a drifting Kรคhler transformation; therefore the real metric does indeed tend to a static solution, and disaster is avoided. However, since we would rather have a fixed endpoint to our Gauss-Seidel method, we โcureโ this drifting due to discretization error by solving the Monge-Ampรจre equation (13), and dynamically determine $`\stackrel{~}{\lambda }`$ by averaging $`detg_{i\overline{ศท}}`$ as we perform the Gauss-Seidel iterations. This procedure then yields a static end solution, and the value $`\stackrel{~}{\lambda }`$ should approach the true analytic $`\lambda `$ as the continuum is approached by increasing the lattice resolution.
It is important to check that the solution to the lattice version of the Monge-Ampรจre equation converges to the continuum solution as the lattice resolution is improved. For this purpose it is useful to have some quantities that can be calculated from the lattice solution, and for which the exact, continuum value is also known. Here we will mention three such quantities. The first is that mentioned directly above, namely the total volume of the manifold, in other words how close $`\stackrel{~}{\lambda }`$ is to $`\lambda `$. The agreement between the analytic and numerical values tests not only the quality of the solution, but also the error in fixing the Kรคhler class. A second quantity that is known exactly is the Euler number $`\chi `$ of the manifold (which of course does not depend on the Kรคhler class). This can be calculated numerically in terms of the Kรคhler potential via a Gauss-Bonnet theorem. Finally, there is the difference between Kรคhler potentials $`K_\alpha K_\beta `$ on the overlap $`U_\alpha U_\beta `$. In subsection 2.1 it was argued that this difference should be set equal to the fixed Kรคhler transformation $`u_{\alpha \beta }`$, as a boundary condition for the $`K_\alpha `$โs. More precisely, however, the boundary condition is imposed only on the *edge* of each patch. For the continuum equation, this is enough to guarantee $`K_\alpha K_\beta =u_{\alpha \beta }`$ also in the interior of $`U_\alpha U_\beta `$, by the uniqueness of the solution to the Monge-Ampรจre equation. The lattice will introduce an error into this equality; conversely, that error measures how well the lattice solution approximates the continuum one.
## 3 Application to Kummer surfaces
For a first application of the techniques described in the previous section, we turned to the lowest-dimensional Calabi-Yau manifold, the K3 surface. K3 has played important roles in algebraic and differential geometry, as well as in string theory. (The lecture notes provide an excellent review of the role of K3 in string theory. See also the more mathematically-oriented review .) The moduli space of Ricci-flat metrics on K3 is 58-dimensional: 40 complex structure moduli and 20 Kรคhler moduli, minus 2 for the hyperkรคhler identifications. One of the moduli is the overall volume of the manifold; the other 57 are non-trivial in the sense that they do not act on the Ricci-flat metric by any straightforward transformation. We studied a particular one-parameter family of K3โs which admit a simple construction and have a high degree of discrete symmetry, which serves to reduce the number of lattice points to be simulated at any given resolution. In the first subsection below, we explain the construction of these K3โs. In the next subsection we describe the results obtained in our simulations, giving several examples of the kind of concrete geometrical information that is available from the explicit form of the metric.
### 3.1 Construction
Among the simplest K3 surfaces to describe explicitly are the so-called Kummer surfaces, which are the orbifold $`T^4/๐_2`$ with its 16 singular points blown up. After explaining this construction, we will specialize to the ones with the largest discrete symmetry group, namely those constructed from a cubical $`T^4`$ with all 16 singular points blown up to the same size. The only free parameters are the size of the $`T^4`$ and the size of the blow-up. We will see that the Kรคhler cone defines a finite range of values for their ratio.
As a complex manifold, the torus $`T^4`$ is parametrized by a pair of complex coordinates $`(z^1,z^2)๐^2`$, identified under translations by a set of four linearly independent vectors $`(v_a^1,v_a^2)๐^2`$ ($`a=1,\mathrm{},4`$),
$$(z^1,z^2)(z^1,z^2)+(v_a^1,v_a^2).$$
(15)
Choosing the vectors $`(v_a^1,v_a^2)`$ fixes the complex structure of the torus (there are equivalences; the complex structure moduli space is actually only 8 real dimensional). The parity map
$$(z^1,z^2)(z^1,z^2)$$
(16)
is compatible with the equivalences (15), so we can quotient the torus by it. Furthermore, it acts holomorphically, so the result, $`T^4/๐_2`$, is a complex manifold. More correctly, itโs an orbifold, since the parity map has fixed points. There are 16 of them, located at $`(z^1,z^2)=\frac{1}{2}_an^a(v_a^1,v_a^2)`$ with $`n^a=0`$ or 1, and each carries an $`A_1`$ (or $`๐^2/๐_2`$) type singularity.
We obtain a smooth complex manifold, known as a Kummer surface, by blowing up the fixed points. Consider for example the fixed point at the origin. To blow it up, remove that point from the manifold and add two new patches, with coordinates $`(y,w)`$ and $`(y^{},w^{})`$ respectively, and transition functions (which are of course holomorphic)
$`(y,w)`$ $`=`$ $`({\displaystyle \frac{1}{y^{}}},w^{}y^2)=({\displaystyle \frac{z^1}{z^2}},{\displaystyle \frac{1}{2}}(z^2)^2),`$ (17)
$`(y^{},w^{})`$ $`=`$ $`({\displaystyle \frac{1}{y}},wy^2)=({\displaystyle \frac{z^2}{z^1}},{\displaystyle \frac{1}{2}}(z^1)^2),`$ (18)
$`(z^1,z^2)`$ $`=`$ $`\pm \sqrt{2w}(y,1)=\pm \sqrt{2w^{}}(1,y^{}).`$ (19)
To avoid complications, the ranges of the new coordinates should be bounded in such a way that they do not include the other fixed points. Each of the 16 fixed points of the orbifold is given its own $`(y,w)`$ and $`(y^{},w^{})`$ patches, with the same transition functions (1719) except that $`(z^1,z^2)`$ is replaced by $`(z^1,z^2)\frac{1}{2}_an^a(v_a^1,v_a^2)`$. Hence we have a total of 33 patches. There are three important points to note about the new $`(y,w)`$ and $`(y^{},w^{})`$ coordinate systems. First, the identification under the orbifold action (16) is automatic in them. Second, the origin has been replaced by the surface $`w=w^{}=0`$, parametrized by $`y=1/y^{}`$. This is a $`๐P^1`$ (or $`S^2`$), and is homologically non-trivial; it is called the exceptional divisor. Finally, the transition functions (1719) all have unit Jacobian. Hence it is quite easy to write down the holomorphic $`(2,0)`$ form:
$$\mathrm{\Omega }=dz^1dz^2=dydw=dw^{}dy^{}.$$
(20)
Kummer surfaces have 8 complex structure moduli (inherited from the $`T^4`$) and 20 Kรคhler moduli (the 4 inherited from the $`T^4`$, plus the size of each of the 16 exceptional divisors). It can be shown that they are special cases of K3 surfaces. (The missing 32 complex structure moduli are due to the fact that we blew up, rather than deformed, the orbifold fixed points.)
Both for simplicity and in order to reduce the number of lattice points simulated, it was advantageous for us to choose highly symmetrical Kummer surfaces. The $`T^4`$ was taken to be cubical; in other words, the periodicities (15) were given by
$$z^1z^1+1z^1+i,z^2z^2+1z^2+i,$$
(21)
while the Kรคhler transformations are those obtained for the flat metric on a cubical $`T^4`$ of side length $`b`$:
$`K(z^1+1,z^2)K(z^1,z^2)`$ $`=`$ $`b^2\left(\mathrm{Re}z^1+{\displaystyle \frac{1}{2}}\right),`$ (22)
$`K(z^1+i,z^2)K(z^1,z^2)`$ $`=`$ $`b^2\left(\mathrm{Im}z^1+{\displaystyle \frac{1}{2}}\right),`$ (23)
$`K(z^1,z^2+1)K(z^1,z^2)`$ $`=`$ $`b^2\left(\mathrm{Re}z^2+{\displaystyle \frac{1}{2}}\right),`$ (24)
$`K(z^1,z^2+i)K(z^1,z^2)`$ $`=`$ $`b^2\left(\mathrm{Im}z^2+{\displaystyle \frac{1}{2}}\right).`$ (25)
The coefficients on the four right-hand sides correspond to the four Kรคhler parameters of $`T^4`$; here since the $`T^4`$ is cubical, they are set equal to a common constant $`b^2`$. Each blown up fixed point has only one Kรคhler modulus; without loss of generality the Kรคhler transformations may be taken as follows:
$`K_{(z^1,z^2)}K_{(y,w)}`$ $`=`$ $`a^2\mathrm{ln}|z^2|,`$ (26)
$`K_{(z^1,z^2)}K_{(y^{},w^{})}`$ $`=`$ $`a^2\mathrm{ln}|z^1|,`$ (27)
$`K_{(y,w)}K_{(y^{},w^{})}`$ $`=`$ $`a^2\mathrm{ln}|y|.`$ (28)
All 16 fixed points were blown up to the same value of the modulus $`a^2`$.
How much discrete symmetry do these Kummer surfaces possess? Our choice of complex structure admits a holomorphic diffeomorphism group of order $`2^8`$, all of which is respected by our choice of Kรคhler class. The generators of this group are as follows: the translations by the vectors $`(\frac{1}{2},0)`$, $`(\frac{i}{2},0)`$, $`(0,\frac{1}{2})`$, and $`(0,\frac{i}{2})`$, which generate a $`๐_2^4`$ group that maps the origin to each of the other fixed points; the rotation $`z^1iz^1`$, which generates a $`๐_4`$ group; the diagonal rotation $`(z^1,z^2)(iz^1,iz^2)`$, which (in view of the identification under parity) generates a $`๐_2`$ group; and the exchange $`(z^1,z^2)(z^2,z^1)`$, which generates another $`๐_2`$. Of course, these generators donโt commute with each other, so the full holomorphic diffeomorphism group is complicated. In addition, the anti-holomorphic diffeomorphism $`(z^1,z^2)(\overline{z}^1,\overline{z}^2)`$ is a symmetry of the real metric. Uniqueness of the solution to the Monge-Ampรจre equation guarantees that every symmetry of the Kรคhler class is an isometry of the Ricci-flat metric in that class. Therefore it is sufficient for us to simulate only the fundamental domain of the symmetry group, a reduction potentially by a factor of $`2^9`$.
As explained in Appendix B, the volume is easily calculated from the intersection matrix of the second de Rham cohomology group $`H_2(๐)`$:
$$V=\frac{1}{2}b^44\pi ^2a^4.$$
(29)
As discussed in subsection 2.3, the volume calculated numerically may be compared against this analytic result to give a measure of how accurately we are fixing our Kรคhler class. More details are given in Appendix A.
Equation (29) shows that, as the size of the blow-up increases, it eats away at the volume of the manifold. This indicates an upper bound $`a^2/b^2<(8\pi ^2)^{1/2}`$. In fact, the real bound is slightly lower: $`a^2/b^2<(4\pi )^1`$. We first noticed this as an empirical fact in our numerical trials, but the reason is not hard to understand. The volume of a holomorphic submanifold depends only on the Kรคhler class, not the metric (the total volume being a special case of this). Therefore, a necessary condition for being inside the Kรคhler cone is that all the holomorphic submanifolds have positive area. Our symmetric Kummer surfaces have three types of holomorphic curves. The first type are the 16 exceptional divisors. From the Kรคhler transformation (28) restricted to the $`w=0`$ surface, one may show that each has area $`A=\pi a^2`$. We thus have the condition $`a^2>0`$. Another holomorphic submanifold is the curve $`\{z^1=C\}\{z^1=C\}`$, which represents the same two-cycle for all values of $`C`$ other than the special values $`C=0,\frac{1}{2},\frac{i}{2},\frac{1}{2}+\frac{i}{2}`$. Using (24,25), its area is $`b^2`$, so we have $`b^2>0`$. In the special cases $`C=0,\frac{1}{2},\frac{i}{2},\frac{1}{2}+\frac{i}{2}`$ the curve passes through points which are outside of the $`z`$ patch, so one has to be more careful. As we show in Appendix B, these represent 4 different two-cycles; each has area
$$\widehat{A}=\frac{1}{2}b^22\pi a^2,$$
(30)
accounting for the above upper limit on $`a^2/b^2`$. Obviously, the same results apply to the curves of constant $`z^2`$, so there are a total of 8 curves of this type. In fact they are rational curves, since they have topology $`T^2/๐_2S^2`$. Each of them intersects each exceptional divisor at a single point, so what is happening as they shrink to zero size is that the exceptional divisors are so large that they touch each other and therefore canโt be blown up any larger.
An isolated rational curve shrinking to zero size implies the formation of an $`A_1`$ orbifold singularity. In the next subsection we will explicitly confirm the formation of these orbifold singularities from our numerical solutions. The simultaneous formation of 8 $`A_1`$ singularities naturally leads to the idea that the manifold in this limit is globally a $`๐_2`$ orbifold. By counting the Euler number, one finds that it must be an orbifold of another, smooth K3. Indeed, it can be shown that it is a rather well-known K3, namely the Fermat quartic surface $`x_1^4+x_2^4+x_3^4+x_4^4=0`$ in $`๐P^3`$, in the Kรคhler class induced from the Fubini-Study metric on $`๐P^3`$.<sup>3</sup><sup>3</sup>3It was shown in that (the blow-up of) the orbifold of the Fermat surface by the $`๐_2`$ action $`(x_1,x_2,x_3,x_4)(x_1,x_2,x_3,x_4)`$ has the same complex structure as the Kummer surfaces we consider. This action has 8 fixed points, and the resulting 8 exceptional divisors are identified with our โshrinking rational curvesโ. It can furthermore be shown that the Kรคhler class on the Kummer surface at $`a^2/b^2=(4\pi )^1`$ lifts to the Kรคhler class on the Fermat surface induced from the Fubini-Study metric on $`๐P^3`$ . Amusingly, it has also been shown that the sigma model with this Kummer surface (or equivalently the orbifold of the Fermat surface) as its target space, at a particular volume and equipped with a particular $`B`$-field, is dual to one whose target space is the Kummer surface at the opposite edge of the Kรคhler cone, $`a^2/b^2=0`$! We are grateful to K. Wendland for very helpful discussions on these issues. Hence our one-parameter moduli space interpolates between the orbifold of $`T^4`$ at $`a^2/b^2=0`$ and the orbifold of the Fermat quartic at $`a^2/b^2=(4\pi )^1`$. For our fixed complex structure, these endpoints of our modulus $`a^2/b^2`$ represent the edges of the Kรคhler cone.
It is worth remarking that an approximate analytical solution to the Einstein equation can be constructed in the limit $`a^2/b^21`$ by smoothly joining the Eguchi-Hanson metric (a Ricci-flat and asymptotically flat metric on the blow-up of $`๐^2/๐_2`$) onto the flat torus metric . As we will see in the next subsection, the numerical method does not perform as well in this regime of parameters because the manifold contains a region of high curvature, namely the vicinity of the exceptional divisors. Thus the numerical and analytic approaches are effective in complementary regimes.
### 3.2 Results
We applied the methods described in Section 2 to find the Ricci-flat metrics on the symmetrical Kummer surfaces constructed in the previous subsection, at 9 different values of the modulus $`a^2/b^2`$. In terms of the combination
$$\alpha =4\pi \frac{a^2}{b^2},$$
(31)
which ranges from 0 to 1, the points for which the metrics were calculated were
$$\alpha =0.03,0.13,0.28,0.50,0.61,0.72,0.79,0.85,0.92.$$
(32)
(In figures 1, 2, 3, and 4 we have left out $`\alpha =0.79`$ for typesetting elegance.) Without loss of generality, the volume of the manifold was fixed to be 1. The computational aspects of the problemโlattice discretization, convergence, etc.โare discussed in detail in Appendix A. Let us simply mention here that, at each of the above values of $`\alpha `$, the metric was computed at four different lattice resolutions, labelled A, B, C, D; B has twice the linear resolution (or 16 times the total number of points) of A, and so on.
In this subsection, we will illustrate the kind of concrete geometric information that is available once the Ricci-flat metric is known. We will focus on three ways to characterize the geometry: the distribution of the Euler density; the induced geometry on the exceptional divisors and other rational curves; and the low-lying spectrum of the Laplacian. If unspecified, results presented are generated from the highest resolution metric, D.
On a Ricci-flat manifold, the simplest non-trivial curvature invariant one can construct is the square of the Riemann tensor (sometimes referred to as the Kretschmann invariant). In four dimensions, this happens to be proportional to the Euler density $`\rho `$:
$$\chi =\sqrt{g}\rho ,\rho =\frac{1}{8\pi ^2}R_{i\overline{ศท}k\overline{l}}R^{i\overline{ศท}k\overline{l}}.$$
(33)
Thus its integral over the manifold is fixed, in this case at 24. Figures 1 and 2 show surfaces of constant $`\rho `$ on the three-dimensional slice $`\mathrm{Im}z^2=0`$ at eight different values of $`\alpha `$.<sup>4</sup><sup>4</sup>4Animations of these plots as a function of $`\mathrm{Im}z^2`$, showing the entire four-dimensional geometry, are available for download at http://schwinger.harvard.edu/~wiseman/k3/. At the smallest value ($`\alpha =0.03`$), the curvature is highly concentrated near the fixed point, and is spherically distributed; here the metric closely approximates the Eguchi-Hanson metric (for which the isosurfaces of $`\rho `$ are spherical in these coordinates, although the full geometry is only axisymmetric). As we increase $`\alpha `$, the curvature spreads out from the fixed point. However, it does not diffuse uniformly over the manifold. Instead, starting in the $`\alpha =0.72`$ figure (third to last), it gathers along the $`z^1=0`$ and $`z^2=0`$ planes. These are the rational curves, discussed in the last subsection (and referred to as $`\widehat{c}_1`$ and $`\widehat{c}_2`$ in Appendix B), that shrink to zero size as $`\alpha `$ approaches 1. By symmetry, the curvature also accumulates at the 6 other shrinking rational curves, $`\{z^1=C\}`$ and $`\{z^2=C\}`$ for $`C=\frac{1}{2},\frac{i}{2},\frac{1}{2}+\frac{i}{2}`$.
In the $`z^i`$ coordinate system used in figures 1 and 2, the exceptional divisor (of the original orbifold) is represented as a single point, namely the origin. This is misleading, since in fact it is topologically an $`S^2`$. Itโs interesting to study how its geometry changes as we vary $`\alpha `$. The induced metric on the exceptional divisor of the Eguchi-Hanson geometry is that of a round sphere. We therefore expect the same to be true for the exceptional divisor of the Kummer surface at small $`\alpha `$. In figure 3, we plot the Ricci scalar of the induced metric on the exceptional divisor. For small values of $`\alpha `$, it is indeed uniform. However, as $`\alpha `$ increases, it becomes non-uniform (its integral is of course fixed at $`8\pi `$): the sphere is becoming prolate. The poles, where the curvature is highest, are at $`y=0`$ and $`y^{}=0`$, in other words where the exceptional divisor intersects the planes $`z^1=0`$ and $`z^2=0`$ respectively; these are the points that are closest to neighboring exceptional divisors.
Similarly, it is interesting to study the induced geometry of the shrinking rational curves, for example the curve $`\{z^1=0\}`$, as shown in figure 4. At $`\alpha =0`$ its geometry is that of a flat $`T^2`$ orbifolded by $`๐_2`$, which is topologically $`S^2`$ but with all the curvature concentrated at the 4 fixed points $`z^2=0,\frac{1}{2},\frac{i}{2},\frac{1}{2}+\frac{i}{2}`$ (picture a square envelope). As $`\alpha `$ increases, the curvature spreads out around the sphere, which becomes rounder and rounder. In the limit $`\alpha 1`$, the geometry in the vicinity of the shrinking rational curve should approach that of an Eguchi-Hanson metric whose exceptional divisor has area $`\widehat{A}=\frac{1}{2}b^22\pi a^2`$. Indeed, we see that for values of $`\alpha `$ approaching 1, the sphere becomes almost completely round. As another test that the geometry is approaching Eguchi-Hanson, in figure 4 we plot the maximum and minimum values of $`\rho `$ on the curve against $`\alpha `$. On an Eguchi-Hanson, the Euler density is constant on the exceptional divisor; the solid curve in that plot is the value of $`\rho `$ on the exceptional divisor of an Eguchi-Hanson of the appropriate size. One sees that as $`\alpha `$ approaches 1 the maximum and minimum values of $`\rho `$ approach each other and the Eguchi-Hanson value.
Given the Ricci-flat metric, one can compute the spectrum of various geometric operators of physical interest, the simplest example being the scalar Laplacian. As a proof of principle, we used our numerical metrics to compute a low-lying eigenvalue (and eigenfunction) of it. The eigenfunctions of the Laplacian on our symmetric Kummer surface can be classified by their eigenvalues under the $`๐_2^4`$ translation subgroup of its full symmetry group. We calculated the lowest eigenvalue in the sector with eigenvalue $`1`$ under all 4 translations. The results are shown in figure 5. The eigenvalue at the orbifold point $`\alpha =0`$, which can be computed exactly, is also shown for comparison. The fact that the eigenvalue increases with $`\alpha `$ can be understood heuristically (at least for small values of $`\alpha `$) in the following way. As the rational curves discussed above shrink, the geodesic distance between a point on the exceptional divisor, such as $`y=w=0`$, and its image under each of the discrete translations decreases, requiring steeper gradients in the eigenfunction. Finally, we would like to point out that, as we are studying a long wavelength eigenfunction, even the lowest resolution (A) computes the eigenvalues to typically within a few percent of the continuum extrapolated value.
We close this section on a technical note. Careful inspection of the isosurfaces plotted for the smallest value of $`\alpha `$ in figure 1 shows a localized deviation from sphericity, caused by small errors near patch overlaps. Similarly, in figure 3 the smallest $`\alpha `$ sphere is slightly less uniform than for the next larger value of $`\alpha `$; in figure 4 for the largest value of $`\alpha `$ we can just see some โringsโ where small errors are introduced due to coordinate patch overlaps; and the distance from the extrapolated continuum values in figure 6 increases near $`\alpha =0`$ or $`1`$. Numerical discretization errors are to be expected, and will be larger where there are higher curvatures. Our solutions containing the highest curvature regions occur near the two orbifold points $`\alpha =0,1`$, explaining why we see the effects mentioned above. Hence, at a fixed resolution, the global quality of the solutions will not be as high near these orbifold points, although the local geometry away from the regions of high curvature should be quite acceptable. Of course, increasing the resolution, the quality of the solution will improve, independently of where we are in moduli space. These issues are discussed in more detail in appendix A.3.
## 4 Discussion
### 4.1 Generalizations
We opened this paper with the question of whether it is possible, in practice, to solve the Einstein equation numerically on a Calabi-Yau manifold. We have shown that the answer is affirmative when the Calabi-Yau is a K3 surface with a high degree of discrete symmetry. In this subsection we will investigate the possibility of generalizing this accomplishment, first to more general K3 surfaces, then to Calabi-Yau three-folds, and finally to Kรคhler manifolds with cosmological constant or with matter.
It is clear that in principle our method extends to a general blow-up of a Kummer surface, and more generally to any K3 surface. It is merely necessary to implement the topology and complex structure with complex coordinates defined appropriately on an atlas of chartsโfor example one could use patches derived from any algebraic construction of K3. However, the question remains as to what resolution is attainable in the absence of large amounts of discrete symmetry. Our highest resolution, $`D`$, simulated approximately $`2\times 10^780^4`$ points. However, due to the high degree of discrete symmetry our one parameter family of K3โs enjoy, every computed point actually represents $`2^9`$ points in the true K3 geometry, and hence we have described the full K3 with an *effective* resolution of around $`400^4`$.
On a current high-end desktop computer (with 1 to 2 gigabytes of memory) one could comfortably increase the total number of points simulated to $`10^8`$. For a K3 with no discrete symmetry this would yield a resolution of $`100^4`$ for the full K3 geometry. To estimate how accuately this could represent the metric, we may compare it to our intermediate resolution, B. With a linear resolution 4 times lower than that for D, the effective resolution for the full geometry is also about $`100^4`$. As seen in subsection 3.2 and appendix A we find that resolution B adequately reproduces the geometry and derived properties, such as the low wavelength eigenmodes of the Laplace operator, provided one is not too close to the edge of the Kรคhler cone. For example, in figure 6 one sees that run B computed the Laplacian eigenvalue with an accuracy of around 1% compared to the extrapolated continuum value. Near the edge of the Kรคhler cone, where regions of high curvature develop in the manifold, the best strategy may be to combine numerical with analytic techniques. For example, at values of the moduli near an orbifold point, one could patch an analytic Ricci-flat metric, such as the Eguchi-Hanson metric, into the region where the high curvature is developing.
If we now wish to move to a Calabi-Yau three-fold, $`10^8`$ points translates to a mere 20 points linearly in each direction. This is over a factor of 2 less in linear resolution than the lowest effective resolution used in this work (run A, which had an effective resolution of $`50^4`$ for the full geometry). This might be acceptable if one were well away from the Kรคhler cone edge, but is certainly rather low. On the other hand, if one were to consider a highly symmetric three-fold, such as a symmetric blow-up of $`T^6/๐_3`$, then one would again expect to attain an effective resolution of around $`100`$ points linearly in the full geometry, and thus again expect accuracy comparable to, or better than, our resolution B.
Moving to the general three-fold appears to be a very challenging task. As we discuss in Appendix A, processing time actually scales rather well with increasing dimension. Instead the problem is limited by storage. In six real dimensions any appreciable increase in linear resolution is extremely costly. Thus whilst $`20^6`$ points would be possible on a desktop computer, $`40^6`$ requires 64 times more memory and is already beyond the abilities of modest clusters. Often a tough computational problem becomes easy in time, as computer memory and speed have closely followed Mooreโs prediction of a doubling every 2 years. However, even assuming that Mooreโs law continues to hold, it will require 12 years to increase the linear resolution by a factor of 2 for the three-folds. It is therefore clear that in order to tackle the general three-fold, one must employ considerably more sophisticated discretization schemes than we have used here, presumably adapting the lattice points to regions of high curvature. Whilst adaptive grids are difficult to implement in the elliptic context, one can easily use fixed grids that increase resolution in areas where curvature is expected to be high.
To summarize, we expect that for a general K3 surface one can obtain very satisfactory results provided one does not wish to probe too near the edge of the Kรคhler cone. For a highly symmetric three-fold similarly high quality results can be expected. However, moving to the general three-fold appears tough and new techniques must certainly be employed.
Having considered the problem of constructing K3โs at arbitrary points in moduli space, we should point out the enormity of the moduli space itself: with 57 directions to explore it is highly implausible that one could ever map the entire space of K3โs. On the other hand it is unlikely that one would ever need to map the entire space. One can imagine wishing to find K3 surfaces with specific properties; presuming the observables of interest vary smoothly over the moduli space, it is plausible that one could scan the moduli space for examples that fit the specific requirements.
Whilst Calabi-Yauโs have large numbers of moduli, spaces with Ricci curvature tend to have fewer of them. An example of this are the four (real) dimensional del Pezzo surfaces dP<sub>n</sub> ($`n=1,\mathrm{},8`$), compact Kรคhler spaces that can be constructed by blowing up $`n`$ points in $`๐P^2`$. Del Pezzos have positive first Chern class, and those with $`n3`$ admit Kรคhler-Einstein metrics with no Kรคhler moduli. The cases of dP<sub>3</sub> and dP<sub>4</sub> are particularly interesting because they also have no complex structure moduli. The Einstein equation can again be reduced to a Monge-Ampรจre equation for the Kรคhler potential. As in our K3 example, two patches would be required to cover each of the $`n`$ blown up points, and three to cover the ambient $`๐P^2`$. Hence we expect the methods we have applied here for K3 to be applicable, and the Kรคhler-Einstein metrics to be attainable to high resolution on a desktop computer (certainly comparable to or better than our resolution C, depending on how many points are blown up in $`๐P^2`$). This would be very satisfying for dP<sub>3</sub> and dP<sub>4</sub> as, with no moduli, one would in principle have constructed these geometries explicitly and completely.
Other physically interesting geometries that are related to Calabi-Yau manifolds are supersymmetric flux compactifications in string theory. A class of type IIB solutions can be constructed as warped products of flat four-dimensional Minkowski spacetime and a Ricci-flat Calabi-Yau three-fold, where the warp factor satisfies a Poisson equation on the Calabi-Yau sourced by fluxes, D-branes, and orientifold planes (see e.g. ). As we discussed above, finding the metric on a generic three-fold is probably out of reach using the methods of this paper. However, considering fluxes on K3 or a highly symmetric three-fold would likely be a manageable task. Having found the Ricci-flat metric, solving the Poisson equation on this geometry is quite simpleโindeed even easier than finding eigenfunctions of the Laplacian as we did earlier in this paper. Studying the solutions to the Poisson equation on the Calabi-Yau background would provide a detailed understanding of how the fluxes backreact on the vacuum geometry, and in particular of how fluxes on adjacent cycles interact.
### 4.2 Lessons for solving general Euclidean geometries
The key simplification in our work has been that of Kรคhler geometry. What are the prospects for constructing general Euclidean geometries? Without Kรคhlerity we require many metric components to describe the geometry, and the Einstein equation becomes complicated. Whilst this may be technically complicated, in principal one would hope that using a harmonic gauge condition would allow the system to be locally solved as an elliptic relaxation problem. Note that one could also use a gauge fixed Ricci flow, but as discussed eariler, it is more efficient to solve the elliptic Ricci flatness condition directly rather than to construct an entire flow when only the endpoint is required. In our K3 example, the complex coordinates on the coordinate patches provide exactly such a local harmonic set of coordinates. However the most challenging aspect of the problem is to understand the global issues, such as residual coordinate freedom, adaptedness of the coordinates, and moduli of the solutions.
Let us for a moment consider the problem of finding the Ricci flat metric of symmetric K3โs as we have done, but ignoring the Kรคhler structure and using only real geometry. We might hope to find harmonic coordinates on the various patches that make up the topology. This would require us to solve the harmonic gauge condition (essentially locally solving Laplace equations) at the same time as the Einstein equation. Presumably this full system is globally elliptic in the case of K3, although for more general geometries we should note that negative modes of the Lichnerowicz operator will exist, as occurs for the Euclidean Schwarzschild solution, and it is unclear how these would affect the situation.
Given the link between the complex coordinates of Kรคhler geometry and the harmonic coordinates natural for Euclidean real geometry, we can make various speculations. We saw that our complex coordinates were well adapted to the symmetric K3 geometry, and one might hope the same to be true for more general harmonic coordinates. As explained above, for the complex coordinates there are no residual holomorphic coordinate transformations; in real geometry, choosing harmonic coordinates one again expects only finitely many residual coordinate freedoms on a compact manifold. Assuming this, one might conclude from our work that we may see the complex structure moduli of the K3 arise simply from the global data required to specify the harmonic coordinates, just as we have fixed the complex structure moduli by taking particular complex coordinates on the manifold. Then for more general compact manifolds, one might plausibly associate physical moduli to global choices when constructing the harmonic coordinates (that are not simply one of the finite residual coordinate transformations).
Clearly about any Ricci flat real geometry one can always linearize metric fluctuations and then, in principle, directly determine the zero modes of the resulting Lichnerowicz operator, and hence determine all physical moduli of the solution. However, this is obviously very complicated to imagine doing in practice, and what we really wish to find is a way to include the moduli as boundary conditions in the elliptic problem as we have done in our Kรคhler example. Assuming our presumptions above about the complex structure moduli hold, the key remaining question is how to understand the Kรคhler moduli of K3 as boundary conditions, without actually making use of the Kรคhler structure. It might be possible to understand this in terms of the volumes of minimal representatives of the two-cycles. However, whilst for K3 this approach would work, for the Kรคhler-Einstein case of the del Pezzos it cannot, since the two-cycles present are not associated with any moduli; hence this approach would not work in general. Thus, finding new ways to understand the Kรคhler moduli and how to actually implement fixing them whilst solving the real geometry Einstein equations for K3, look to be important questions. If they can be addressed, it might allow one to understand the moduli of general real geometries, and enable explicit metrics to be found in very general Euclidean geometry-matter systems.
### 4.3 Applications
We now briefly discuss possible applications for the numerical construction of geometries. We consider first mathematical, then formal physical, and finally phenomenological applications.
There are various outstanding mathematical conjectures regarding the geometry of Calabi-Yau manifolds. Most notorious is the Strominger-Yau-Zaslow conjecture, that Calabi-Yaus with mirrors can be constructed as toric fibrations and that mirror symmetry acts by T-duality on the fibers . A related conjecture concerns mean curvature flows and special Lagrangian submanifolds in Calabi-Yau manifolds . In principle these conjectures can be tested directly on any Calabi-Yau for which the Ricci-flat metric is known explicitly.
From the point of view of physics, explicit constructions of Ricci-flat Calabi-Yau metrics may help us learn more about the sigma model description of geometry. The Ricci-flat metric on K3 is the target space of an $`๐ฉ=(4,4)`$ non-linear sigma model. Due to the high degree of supersymmetry the classical Ricci flat metric receives no perturbative or non-perturbative corrections in $`\alpha ^{}`$ . Thus the geometries we have constructed in this paper can, remarkably, be viewed as fully quantum geometries from the viewpoint of this sigma model. Knowing these metrics then in principle allows one to compute properties of the quantum sigma modelโfor example, the spectra of operators on the target manifold correspond to conformal weights on the worldsheet.
The Ricci-flat Calabi-Yau three-folds are again target spaces of sigma models. However, now the classical geometry only gives the leading $`\alpha ^{}`$, or large volume, approximation to the true quantum geometry. Understanding how $`\alpha ^{}`$ corrections modify the classical geometry is an important physical issue. Supersymmetry implies that these corrections preserve the Kรคhlerity of the metric, and hence they will appear as higher derivative terms modifying the Monge-Ampรจre equation for the Kรคhler potential. Each higher derivative term will make this equation less local, but in principle we may still apply the same local iterative methods we have used here to solve it. (Presumably on a compact manifold, one could in principle include infinitely high derivative terms if their form were known.) From a discretized viewpoint, if we linearized the equation about some background, we would find an $`N\times N`$ operator ($`N`$ being the total number of lattice points) and the structure of its matrix representation will no longer be sparse. However, the terms that fill in the zero components in the sparse Monge-Ampรจre case will be small, being down by factors of $`\alpha ^{}`$. Hence our Gauss-Seidel iterative methods may still work, although obviously evaluation of the equation at each point will take much longer.
Finally, from a phenomenological point of view, being able to compute metrics is crucial for actually making contact with low energy physics in string theory. Whilst for simple vacuum Calabi-Yau reductions it is possible to compute the entire low energy effective action using only topological data , as soon as matter is added to the compactification manifold this is no longer true. The simplest example is adding a single brane that is localized in the compact space. The moduli space metric for this brane, and hence the kinetic term for its position in the low-energy action, is simply given by the metric on the internal space . Thus from a physics standpoint, being able to compute the low energy action of geometric reductions is a strong motivation to further understand and improve the numerical geometry methods we are exploring here.
It is worth mentioning that in TeV fundamental scale senarios it is possible that in a few years the LHC might directly probe not just the low energy action, but also high energy excitations on the internal space, such as Kaluza-Klein modes. Although the LHC would likely not measure a sufficient number for any detailed spectroscopy of the internal space, future colliders would then be able to measure these massive excitations accurately. If one could extend our methods to the general three-foldโand surely this would be enough motivation to direct serious resources to itโthen one might scan the moduli space of Calabi-Yaus, presumably with fluxes, to find candidate geometries matching the observed resonances.
###### Acknowledgments.
We would like to thank the following people, who have given us invaluable help in this project: A. Adams, F. Denef, J. Distler, M. Douglas, D. Gaiotto, S. Gukov, J. Hartle, B. Julia, B. Kors, J. Minahan, D. Morrison, L. Motl, W. Nahm, A. Neitzke, C. Nunez, R. Reinbacher, A. Sen, J. Sparks, A. Strominger, P. Tripathy, B. Wecht, K. Wendland, and S.-T. Yau. M.H. is supported by a Pappalardo Fellowship, and by the U.S. Department of Energy through cooperative research agreement DF-FC02-94ER40818. T.W. is supported by the David and Lucile Packard Foundation, grant number 2000-13869A.
## Appendix A Details of numerical construction
### A.1 Construction of the atlas and initial data
As discussed in section 3.1 our one parameter family of K3โs have many discrete symmetries. In particular these imply we may take the Kรคhler potential to be identical in each of the 16 regions describing the blow up of the torus fixed points. These regions are each described in terms of 2 coordinate patches given by $`w,y`$ and $`w^{},y^{}`$ and the symmetry $`z^1z^2`$ implies that the Kรคhler potential may be taken to be identical in each of these. Thus we reduce our problem to one patch describing the fundamental domain of the torus using $`z^1,z^2`$ coordinates, where we orbifold by the identification $`(z^1,z^2)(z^1,z^2)`$, and one patch describing (half of) the blow up of the fixed point contained in that fundamental domain using coordinates $`w,y`$. We term these patches the โTorus patchโ and the โEguchi-Hanson patchโ. To begin describing the geometry of our atlas we define,
$$\sigma =|z^1|^2+|z^2|^2=2|w|\left(1+|y|^2\right)$$
(34)
We take our torus patch to cover the coordinate range of the fundamental domain, but exclude the region near the blown up fixed point so,
$$\frac{1}{4}\mathrm{Re}z^{1,2}\frac{1}{4},\frac{1}{4}\mathrm{Im}z^{1,2}\frac{1}{4}\mathrm{and}\sigma \sigma _{min}$$
(35)
Outside the boundaries of this domain we must act with the translations $`z^{1,2}z^{1,2}\pm (1,i)/2`$ to map the point back into the domain. Note that when we do this, we must also ensure we perform the torus Kรคhler transformation derived from 25 reduced to this domain. The Eguchi-Hanson patch is taken to have coordinate range,
$$0|y|y_{max}\mathrm{and}\sigma \sigma _{max}$$
(36)
In order to cover the manifold we must ensure $`y_{max}1`$, and $`\sigma _{max}\sigma _{min}`$ with $`\sigma _{max}<1/4^2`$ to avoid complicated multiple overlaps.
In order to evaluate our Monge-Ampรจre equation at the edge of a coordinate patch we must necessarily compute derivatives involving Kรคhler potential data from neighbouring coordinate patches. Once we have finite differenced our patches this will require the Kรคhler potential at some point to be computed from the neighbouring patch and then Kรคhler transformed into the original patch. The coordinate location in the neighbouring patch will not necessarily fall on a lattice point and therefore we will need to perform interpolation to compute the desired Kรคhler potential <sup>5</sup><sup>5</sup>5Note that whilst extrapolation is more economical as we would require no patch overlap, it is also potentially dangerous as it is likely to introduce spurious data, and since specifiying the correct data is our key concern we have opted to use interpolation which takes a little more storage (due to the patches overlapping) but removes the risk of specifying data incorrectly.. Thus we require our patches to overlap sufficiently in order to perform our necessary interpolations.
For example, suppose when we evaluate a derivative at the boundary of the Eguchi-Hanson patch we require knowing the Kรคhler potential at a point still with $`\sigma <\sigma _{max}`$ but now $`|y|>y_{max}`$. Then we must use the coordinate patch $`w^{},y^{}`$ \- which due to the discrete symmetries is identical to the $`w,y`$ one. Explicitly, we transform to the $`w^{},y^{}`$ coordinates where still $`\sigma <\sigma _{max}`$, but now $`|y^{}|<y_{max}`$, so the point does indeed lie within this $`w^{},y^{}`$ patch. We find the Kรคhler potential in this patch using the necessary interpolation, and then return to our original coordinate patch $`w,y`$ by performing the necessary Kรคhler transformation 28. Similarly, at the large $`\sigma `$ boundary of the Eguchi-Hanson patch, or small $`\sigma `$ boundary of the torus patch we will find the Kรคhler potential in that patch by interpolating from the other and performing the appropriate Kรคhler transformations 26, 27.
We note that in fact we actually require only one coordinate patch, as the Eguchi-Hanson patch can quite satisfactorily represent the torus region. However, the torus patch boundary conditions, essentially derived from 22-25 become complicated and non-local in the $`w,y`$ coordinates, and we have found it simpler to use the two patches above, rather than the minimal choice of one patch.
Even after the reduction to these 2 patches, and the coordinate domains above, we still have discrete symmetries remaining. The orbifold symmetry, holomophic isometries $`z^jiz^j`$, $`z^1z^2`$ and anti-holomorphic isometry $`z^{1,2}\overline{z}^{1,2}`$ further reduce the torus coordinate domain (in addition to $`\sigma \sigma _{min}`$) to,
$`0\mathrm{Re}z^{1,2}`$ $``$ $`{\displaystyle \frac{1}{4}},0\mathrm{Im}z^{1,2}{\displaystyle \frac{1}{4}}`$
$`\mathrm{and}(z^1,z^2)`$ $``$ $`(z^2,z^1)`$
$`(z^1,z^2)`$ $``$ $`(\overline{z}^1,\overline{z}^2)`$ (37)
In the Eguchi-Hanson patch these isometries allow us to further reduce 36 to,
$$0\mathrm{Re}w,y0\mathrm{Im}w,y$$
(38)
The Kรคhler potential does not transform under these discrete isometries, and therefore if we now require a point to be evaluated outside our reduced coordinate domains we simply use these discrete symmetries to map the point back into our reduced coordinate domain.
In the data we present we have chosen a fixed geometry for the patches and their overlaps. This is independent of the numerical resolution of the discretization so that for convergence testing we are comparing like with like. As discussed in section 2.3 it also allows us to directly compare the Kรคhler potential in the same overlapping regions as we vary the resolution which provides a check of numerical convergence. The parameters we have chosen are,
$$y_{max}=1.25\sigma _{min}=\frac{1}{4^2}\times 0.32\mathrm{and}\sigma _{max}=\frac{1}{4^2}\times 0.60$$
(39)
Before we begin relaxing the Monge-Ampรจre equation we require some smooth initial data compatible with our choice of Kรคhler parameters $`a,b`$ (although in fact we have found that in some cases even taking non-smooth initial data the algorithm still converges). We construct this by taking an initial guess Kรคhler potential to have the behaviour of the Eguchi-Hanson potential near $`\sigma =0`$ in the Eguchi-Hanson patch, and which then interpolates smoothly up to second derivatives to behave as the flat torus potential for $`\sigma >\sigma _{max}`$.
The Eguchi-Hanson potential in the coordinate patch $`w,y`$ is given by,
$$K_{(y,w)}^{\mathrm{EH}}=\frac{1}{2}\sqrt{\sigma ^2+a^4}+\frac{a^2}{2}\mathrm{ln}\frac{\mathrm{ln}\left(1+|y|^2\right)}{1+\sqrt{1+\frac{\sigma ^2}{a^4}}}$$
(40)
and the flat torus potential is simply $`K^{\mathrm{torus}}=\frac{1}{2}\sigma b^2`$. For our interpolation, in our Eguchi-Hansen patch (with $`\sigma <\sigma _{max}`$) we take,
$`K_{(y,w)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}a^2\mathrm{log}(1+|y|^2)+k_0+k_2\sigma ^2+k_4\sigma ^4`$
$`\mathrm{with}k_0`$ $`=`$ $`{\displaystyle \frac{1}{16}}\left(6a^2+3\sigma _{max}b^28a^2\mathrm{log}\sigma _{max}\right)`$
$`k_2`$ $`=`$ $`{\displaystyle \frac{1}{8\sigma _{max}^2}}\left(3\sigma _{max}b^24a^2\right)`$
$`k_4`$ $`=`$ $`{\displaystyle \frac{1}{16\sigma _{max}^4}}\left(\sigma _{max}b^22a^2\right)`$
and in the torus patch with $`\sigma >\sigma _{min}`$ we take,
$`\sigma \sigma _{max}K_{(z1,z2)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}a^2\mathrm{log}\sigma +k_0+k_2\sigma ^2+k_4\sigma ^4`$
$`\sigma >\sigma _{max}K_{(z1,z2)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\sigma b^2`$
The constants above ensure the Kรคhler potential has smooth second derivatives at $`\sigma =\sigma _{max}`$. However, it is easy to show that the above does not define a positive Kรคhler form over the whole range $`\alpha =0`$ to $`1`$. As discussed in the main text, one advantage in solving the Monge-Ampรฉre equation directly, rather than performing Ricci flow, is that the Kรคhler form need not be positive. Hence the simple interpolation above suffices as initial data.
### A.2 Discretization, memory and time requirements
We discretize our system in the most naive way. In both patches we discretize by creating a uniform lattice in the real and imaginary parts of each complex coordinate. We then use second order finite differencing to implement the Monge-Ampรฉre equation and third order accurate interpolation at the patch overlaps.
At each step of our relaxation, we firstly interpolate the values of the Kรคhler potential at the very edges of a patch from the appropriate neighbouring patches, and secondly perform one iteration of the Gauss-Seidel generalizated to our Monge-Ampรจre equation. When updating each point, as a by-product we may quickly compute the value of $`detg`$ at that point. During the Gauss-Siedel update, we keep a running average of this quantity over all lattice points, and then use this as the value of $`\stackrel{~}{\lambda }`$ for the next step. This ensures our Gauss-Seidel iterations asymptote to a fixed solution.
In this paper we have used 4 different resolutions each differing in linear resolution by a factor of 2. In the the torus patch we discretize with equal lattice spacing $`dz`$ in the real and imaginary $`z^1`$ and $`z^2`$ directions. In the Eguchi-Hanson patch we discretize with spacing $`dw`$ in the real/imaginary $`w`$ directions and $`dy`$ in the real/imaginary $`y`$ directions. We label the various resolutions $`AD`$, and they are specified as,
$$\begin{array}{ccccccccc}\hfill Adz& =& 0.025\hfill & \hfill dw& =& 0.003125\hfill & \hfill dy& =& 0.25\hfill \\ \hfill Bdz& =& \frac{1}{2}0.025\hfill & \hfill dw& =& \frac{1}{2}0.003125\hfill & \hfill dy& =& \frac{1}{2}0.25\hfill \\ \hfill Cdz& =& \frac{1}{2^2}0.025\hfill & \hfill dw& =& \frac{1}{2^2}0.003125\hfill & \hfill dy& =& \frac{1}{2^2}0.25\hfill \\ \hfill Ddz& =& \frac{1}{2^3}0.025\hfill & \hfill dw& =& \frac{1}{2^3}0.003125\hfill & \hfill dy& =& \frac{1}{2^3}0.25\hfill \end{array}$$
(43)
Remembering that the coordinates $`z`$ range from $`00.25`$ for our reduced domain in the torus patch, this yields 80 points along a side of the torus in resolution $`D`$.
In the torus patch we implement the 2 identifications in A.1 imperfectly by simply only storing points with $`\mathrm{Im}(z^2)>\mathrm{max}(\mathrm{Re}(z^1),\mathrm{Im}(z^1))`$. This is rather convenient, but doesnโt fully take advantage of these discrete symmetries. Asymptotically this yields a reduction of 3 rather than the optimal value 4 for these 2 identifications. With this minor imperfections in mind, the total number of points stored in computer memory for each resolution to represent the manifold is then,
$$\begin{array}{ccc}\hfill \text{number of points}A& =& 6\times 10^3\hfill \\ \hfill B& =& 7\times 10^4\hfill \\ \hfill C& =& 1\times 10^6\hfill \\ \hfill D& =& 2\times 10^7\hfill \end{array}$$
(44)
The Gauss-Seidel scheme is implemented by solving the discrete Monge-Ampรจre equations at each point in the lattice. We found that under-relaxation was not required for stability and the scheme converged stably. Interestingly we found that we could not over-relax the equation with any appreciable over-relaxation rendering the scheme unstable. Presumably this is an effect associated with the patch boundaries rather than their interiors which we expect to behave in an analogous manner to the Poisson equation. Thus in principle it might be possible to include an over-relaxation parameter that varied over the patch to be one at the edges, but larger than one in the interior.
We measure the distance from convergence by computing the maximum update of the Kรคhler potential in a given Gauss-Seidel iteration over the whole lattice. This is equivalent to computing the maximum violation of the discretized Monge-Ampรจre equation.
When this number falls below $`10^{12}`$ we classify the solution as having relaxed. Certainly for any quantity we have computed, there is no further change if the solutions are subjected to further iterations of the Gauss-Seidel scheme. The time taken for our implementation of the algorithm from initial guess to the relaxed condition is shown in figure 7 for the various resolutions $`AD`$ on a standard desktop computer (3Ghz Pentium, 500 Mb). After a little relaxation, as for the Poisson equation, the number of Gauss-Seidel iterations required to improve the discretized error by a factor of 10 quickly tends to a constant. This is also shown in the same figure.
We see that both these quantities exhibit only a weak dependence on the position in moduli space that we choose. The lowest resolution $`A`$ relaxes in seconds while our highest resolution $`D`$ requires up to a week.
The time taken to relax using the local Gauss-Seidel scheme can be estimated as going as $`N^{1+2/d}`$ in $`d`$ real dimensions where $`N`$ is the total number of lattice points. Every iteration takes a time of order $`N`$, and the total number of steps can be estimated as $`N^{2/d}`$ by considering the spectral radius of the linearized Laplace operator. We see this scaling is certainly consistent with the results of figure 7.
At every step in the iteration we must also perform an interpolation of points at the edge of each coordinate patch from its neighbours. We only require the edge points to be interpolated that are required by our second order differenced Monge-Ampรจre equation. Thus the work involved scales as a codimension one quantity, namely as $`N^{(d1)/d}`$. In practice our third order interpolation actually takes considerable time. Asymptotically at large $`N`$, however, it will obviously become subdominant to the Gauss-Seidel iteration time which scales as $`N`$.
It is very interesting to note that as we move to higher dimensions, the total relaxation time more and more closely approaches $`N`$. Thus the advantage in using highly non-local schemes such as multi-grid to improve convergence times (typically to $`N\mathrm{log}N`$) becomes considerably reduced. In this sense we may claim that the problem of extending our methods to Calabi-Yau 3-folds is storage limited rather than speed limited.
In a memory limited problem, it is rather natural to move to higher order methods. Our implementation uses second order finite differencing on the Monge-Ampรฉre equation. However, we might hope for improved convergence to the continuum if we were to use $`4^{th}`$ order discretization of the Monge-Ampรจre equation (and also for interpolating between patches). We did try this in our case of K3, but the additional time required by each iteration slowed the total convergence time to approximately the same as the next higher resolution using second order differencing. In order to procede to our highest resolution we decided to stay with second order differencing. Bear in mind that whilst a higher order method approaches the continuum more quickly, in order to resolve short length scales one requires high resolutions, so to get accurate results near the orbifold regions of our moduli space, we require the highest resolutions possible.
With a particular physics or maths question in mind, one might attempt computations using more resources and improved stamina than we have used here. Then the fundamental problem of limited storage should probably be tackled by both a combination of more efficient discretization, and also higher order methods. The increased cost in processor time could certainly be ameliorated by some form of parallelisation - which is well suited to this problem which, afterall, naturally divides into coordinate patches.
The eigenfunction of the scalar Laplace operator presented in section 3.2 was computed using a naive iterative scheme. The initial guess for the eigenfunction $`\psi _{(\alpha )}`$ with odd parity under the $`Z_2^4`$ translation isometry ($`z^iz^i\pm 1/2`$) was taken to be that for the torus orbifold, $`\alpha =0`$, where,
$$\psi _{(0)}=\mathrm{cos}2\pi \mathrm{Re}z^1\mathrm{cos}2\pi \mathrm{Im}z^1\mathrm{cos}2\pi \mathrm{Re}z^2\mathrm{cos}2\pi \mathrm{Im}z^2$$
(45)
giving an eigenvalue $`4(2\pi /b)^4`$. The eigenvalue equation was solved using Gauss-Seidel iteration everywhere but at one point, with the eigenvalue being determined dynamically from the condition that the eigenfunction be smooth at that point. Of course, explicit independence of the actual point chosen was checked. This method is very simple to implement, but is really only suited to finding the lowest eigenfunction in a parity sector. More sophisticated (but standard) methods would be required to computer higher eigenfunctions.
### A.3 Convergence tests
We now breifly present data that demonstrates increasing numerical resolution improves observed quantities in a manner consistent with a second order approach to the continuum.
Firstly as discussed in the main text, we compute a numerical $`\stackrel{~}{\lambda }`$ rather than using the analytic value $`\lambda `$. This ensures the numerical Monge-Ampรจre equations converges to a static end point. We may then compare this numerical determination of $`\lambda `$ with the true analytic value. In figure 8 we plot the difference between these values for the various resolutions as a function of position in moduli space. We see the error is greatest near the orbifold point at $`\alpha =1`$ as expected. We clearly see that the errors improve with increasing resolution consistent with second order scaling - we remind the reader that the 4 resolutions $`AD`$ each differ by a linear resolution factor of 2.
In figure 9 we plot the integrated Euler density for the resolutions $`BD`$. The result, the Euler number, has true value 24. We clearly see here that increasing resolution does indeed improve the numerical determination of this quantity as we would hope for. The lowest resolution gives a very poor estimate of the Euler number and we have not included it here. Resolution $`B`$ still provides a rather poor estimate. We see the error in the Euler number is practically quite small for resolution $`C`$ provided we are not too near either orbifold point, and resolution $`D`$ gives errors of less than $`0.1\%`$ over most of the moduli space.
Finally our patches geometrically overlap in coordinate space. Therefore as discussed in the main text, once we have found a solution, we can test its quality by computing the error in the Kรคhler potential in the overlapping regions (obviously taking into account the relevent Kรคhler transformation between the 2 patches). In figure 10 we show the maximum error found by comparing 2 different patch overlaps; the overlap of the Eguchi-Hanson patch with itself, and then its overlap with the torus patch. We see that this maximum error again decreases consistent with second order scaling as resolution is increased.
## Appendix B Homology of Kummer surfaces
In this appendix we derive some properties of the second homology group of Kummer surfaces that were used in the main text. We also show how to relate the natural basis for the homology to the standard basis for the integral homology of K3. As far as we know this explicit relation is new.
K3 has second Betti number $`b_2=22`$. In the Kummer construction, 6 two-cycles are inherited from $`T^4`$, and the other 16 are the exceptional divisors at the blown up fixed points. The second homology of $`T^4`$ is generated by the 2 holomorphic curves $`\{z^1=C\}`$ and $`\{z^2=C\}`$, as well as the 4 non-holomorphic curves $`\{\mathrm{Re}z^1=C^1,\mathrm{Re}z^2=C^2\}`$, $`\{\mathrm{Re}z^1=C^1,\mathrm{Im}z^2=C^2\}`$, $`\{\mathrm{Im}z^1=C^1,\mathrm{Re}z^2=C^2\}`$, and $`\{\mathrm{Im}z^1=C^1,\mathrm{Im}z^2=C^2\}`$. When taking the orbifold, one may choose the constant(s) in a such a way that the curve either passes through or avoids the fixed points (e.g. the holomorphic curves pass through four fixed points if $`C=0`$ but avoids them if $`C=1/4`$). In the latter case, one must include the image under the orbifold, e.g. $`\{z^1=C\}\{z^1=C\}`$, to obtain a two-cycle in the orbifold. (The cycles that pass through the fixed points are linear combinations of those that donโt and the exceptional divisors.) For our purposes, it is simplest to take as a basis the 6 cycles that avoid the fixed points, which we will refer to as โtorus cyclesโ, along with the 16 exceptional divisors. We refer to these cycles as $`c_I`$, where $`c_{1,2}`$ are the holomorphic torus cycles, $`c_{3,4,5,6}`$ are the non-holomorphic torus cycles, and $`c_{7,\mathrm{},22}`$ are the exceptional divisors (which are holomorphic).
It is straightforward to write down the intersection matrix in this basis. The torus cycles intersect each other in pairs, with intersection number 2. For example $`c_1`$ intersects $`c_2`$ at two points, $`(C,C)`$, and $`(C,C)`$ (of course we donโt count $`(C,C)`$ and $`(C,C)`$ separately). By construction, none of the torus cycles intersect the exceptional divisors. The latter do not intersect each other, but have self-intersection $`2`$ (since they are topologically $`๐P^1`$โs). All in all, we find the following block diagonal intersection matrix:
$$h_{IJ}=\mathrm{\#}(c_I,c_J)=2\left[\begin{array}{cccc}U& & & \\ & U& & \\ & & U& \\ & & & I_{16}& \end{array}\right],U=\left[\begin{array}{cc}0& 1\\ 1& 0\end{array}\right].$$
(46)
From the Kรคhler transformations (2228) the periods of the Kรคhler form $`j_I=_{c_I}J`$ may be computed for the symmetric Kummer surfaces:
$$j_I=\{\begin{array}{cc}b^2,\hfill & I=1,2\hfill \\ 0,\hfill & I=3,4,5,6\hfill \\ \pi a^2,\hfill & I=7,\mathrm{},22\hfill \end{array}.$$
(47)
In the case of the holomorphic curves ($`I=1,2,7,\mathrm{},22`$) these periods are their areas. In terms of the periods the volume is
$$V=\frac{1}{2}JJ=\frac{1}{2}h^{IJ}j_Ij_J=\frac{1}{2}b^44\pi ^2a^4,$$
(48)
where $`h^{IJ}`$ is the inverse intersection matrix,
$$h^{IJ}=\frac{1}{2}\left[\begin{array}{cccc}U& & & \\ & U& & \\ & & U& \\ & & & I_{16}& \end{array}\right].$$
(49)
Consider now the holomorphic curve $`\widehat{c}_1=\{z^1=0\}`$. This intersects $`c_2`$ at one point, $`(0,C)`$, but none of the other torus cycles. It also intersects, at one point each, 4 of the exceptional divisors, namely those located at $`(0,0)`$, $`(0,\frac{1}{2})`$, $`(0,\frac{i}{2})`$, $`(0,\frac{1}{2}+\frac{i}{2})`$, which we will call $`c_{7,8,9,10}`$. Using the intersection matrix, we have
$$\widehat{c}_1=\frac{1}{2}\left(c_1c_7c_8c_9c_{10}\right).$$
(50)
Hence the area of this curve is
$$\widehat{A}=_{\widehat{c}_1}J=\frac{1}{2}\left(j_1j_7j_8j_9j_{10}\right)=\frac{1}{2}b^22\pi a^2,$$
(51)
as claimed in subsection 3.1.
Now let us return to the general Kummer surface. The two-cycles $`c_I`$ are obviously integral. However, since the inverse intersection matrix (49) is not composed of integers, they do not form a basis for the integral homology $`H_2(\text{K3},๐)`$, but only a sublattice of it. The standard basis $`\{c_M^{}\}`$ for $`H_2(\text{K3},๐)`$ is defined to have intersection matrix
$$h_{MN}^{}=\mathrm{\#}(c_M^{},c_N^{})=\left[\begin{array}{ccccc}U& & & & \\ & U& & & \\ & & U& & \\ & & & E_8& \\ & & & & E_8\end{array}\right],$$
(52)
where $`E_8`$ represents the Cartan matrix of that group,
$$E_8=\left[\begin{array}{cccccccc}2& 1& & & & & & \\ 1& 2& 1& & & & & \\ & 1& 2& 1& & & & \\ & & 1& 2& 1& & & \\ & & & 1& 2& 1& 1& \\ & & & & 1& 2& & \\ & & & & 1& & 2& 1\\ & & & & & & 1& 2\end{array}\right].$$
(53)
The change of basis relating the $`\{c_I\}`$ to the $`\{c_M^{}\}`$ is as follows:
$$c_M^{}=M_{M}^{}{}_{}{}^{I}c_I,$$
(54)
where
$$\begin{array}{c}2M_{M}^{}{}_{}{}^{I}=\hfill \\ \hfill \left[\begin{array}{cccccccccccccccccccccc}40& 20& 11& 2& 14& 0& 13& 4& 8& 1& 8& 1& 11& 20& 20& 20& 2& 0& 0& 0& 0& 2\\ 46& 23& 11& 2& 15& 0& 14& 4& 10& 0& 9& 1& 13& 23& 22& 24& 2& 0& 0& 0& 0& 2\\ 8& 4& 0& 0& 0& 0& 2& 0& 2& 0& 2& 0& 2& 4& 4& 4& 0& 0& 0& 0& 0& 0\\ 6& 3& 1& 1& 0& 0& 1& 0& 1& 1& 1& 1& 2& 3& 3& 3& 1& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 2& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0\\ 2& 1& 1& 1& 1& 1& 1& 0& 0& 1& 0& 0& 1& 1& 1& 1& 0& 0& 1& 0& 0& 1\\ 4& 1& 0& 0& 1& 0& 1& 1& 0& 0& 0& 0& 1& 1& 2& 2& 0& 0& 0& 0& 0& 0\\ 0& 1& 1& 0& 1& 0& 0& 0& 0& 0& 1& 1& 1& 1& 0& 0& 0& 0& 0& 0& 0& 0\\ 4& 2& 1& 0& 0& 0& 1& 0& 0& 1& 2& 1& 1& 2& 2& 2& 0& 0& 0& 0& 0& 0\\ 2& 1& 1& 0& 1& 0& 0& 0& 1& 1& 0& 0& 1& 1& 1& 1& 0& 0& 0& 0& 1& 1\\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 2& 0\\ 14& 7& 3& 0& 5& 0& 4& 1& 3& 0& 3& 0& 4& 7& 7& 7& 0& 0& 0& 0& 1& 1\\ 22& 11& 5& 0& 5& 0& 6& 1& 5& 0& 5& 0& 6& 11& 11& 11& 0& 0& 0& 0& 1& 1\\ 7& 3& 2& 1& 3& 0& 2& 2& 1& 1& 1& 0& 2& 3& 3& 4& 1& 0& 0& 0& 0& 0\\ 2& 1& 0& 0& 1& 0& 0& 0& 0& 0& 1& 1& 1& 1& 1& 1& 0& 0& 1& 1& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 2& 0& 0\\ 4& 2& 1& 0& 0& 0& 1& 0& 1& 0& 1& 0& 1& 2& 2& 2& 1& 1& 1& 1& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 2& 0& 0& 0& 0\\ 2& 1& 1& 0& 1& 0& 0& 0& 0& 0& 0& 0& 0& 2& 1& 1& 1& 1& 0& 0& 0& 0\\ 2& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 2& 0& 0& 0& 0& 0& 0& 0& 0\\ 2& 2& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 2& 2& 2& 0& 0& 0& 0& 0& 0\\ 15& 7& 4& 1& 5& 0& 5& 2& 3& 0& 3& 1& 5& 7& 7& 7& 0& 0& 1& 0& 0& 1\end{array}\right].\end{array}$$
(55)
(Thus one has $`h_{MN}^{}=M_{M}^{}{}_{}{}^{I}M_{N}^{}{}_{}{}^{J}h_{IJ}`$.) In deriving this change of basis we made use of results in . To our knowledge its explicit form is new.
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# Internal Consistency of Fault-Tolerant Quantum Error Correction in Light of Rigorous Derivations of the Quantum Markovian Limit
## I Introduction
The theory of fault-tolerant quantum error correction (FT-QEC) is one of the pillars that the field of quantum information rests on. Starting with the discovery of quantum error correcting codes Shor:95 ; Steane:96a , and the subsequent introduction of fault tolerance Shor:96 , this theory has been the subject of many improvements and important progress Knill:96 ; Knill:98 ; Aharonov:96b ; Aharonov:99 ; Zalka:96 ; Gottesman:97b ; Gottesman:99 ; Gottesman:99a ; Preskill:97a ; Steane:99a ; Steane:03 ; Knill:04 ; Reichardt:04 ; Knill:05 , leading to the well-known error correction threshold condition. Most recently, work by Steane Steane:03 and Knill Knill:05 (see also Reichardt Reichardt:04 ) has pushed the threshold down to values that are claimed to be very close to being within experimental reach. A notable feature of much of the work on FT-QEC is that the error models are *phenomenological*. By this we mean that the underlying models often do not start from a Hamiltonian, microscopic description of the system-bath interaction, but rather from a higher level, effective description, most notably that of Markovian dynamics. E.g., Knill writes: โWe assume that a gateโs error consists of *random, independent* applications of products of Pauli operators with probabilities determined by the gateโ (our emphasis) Knill:05 . This approach is natural given the considerable difficulty of obtaining error thresholds starting from a purely Hamiltonian description. Nevertheless, Hamiltonian approaches to decoherence management in a fault-tolerant setting have been pursued, e.g., a mixed phenomenological-Hamiltonian treatment of FT-QEC MohseniLidar:05 ; Terhal:04 ; Aliferis:05 ; Aharonov:05 , and a fully Hamiltonian study of fault tolerance in dynamical decoupling KhodjastehLidar:04 . Also noteworthy are recent mixed phenomenological-Hamiltonian *continuous time* treatments of QEC Ahn:01 ; Sarovar:04 ; Sarovar:05 .
Here we are concerned with a critical re-evaluation of the physical assumptions entering the theory of FT-QEC. We scrutinize, in particular, the consistency of the assumption of Markovian dynamics within the larger framework of FT-QEC. We point out that there may be an inherent inconsistency in the theory of Markovian FT-QEC, when viewed from the perspective of the validity of the Markovian approximation. We begin by briefly reviewing, in Section II, a set of minimal and standard, universally agreed upon assumptions made in Markovian FT-QEC theory. We then review, in Section III, the derivation of Markovian master equations, emphasizing the physical assumptions entering the Markovian approximation, in particular the requirement of consistency with thermodynamics. Having delineated the set of assumptions entering FT-QEC and the quantum Markov approximation, we discuss in Section IV the internal consistency of Markovian FT-QEC theory. We point out where according to our analysis there is an inconsistency, and discuss possible objections. In Section VI we then discuss how one may overcome the inconsistency using a variety of alternative approaches, including adiabatic quantum computing (QC), holonomic QC, topological QC, and recent work on FT-QEC in a non-Markovian setting Terhal:04 ; Aliferis:05 ; Aharonov:05 . We conclude in Section VII.
## II Review of Standard Assumptions of FT-QEC
The following are a set of minimal assumptions made in the theory of FT-QEC Shor:96 ; Knill:96 ; Knill:98 ; Aharonov:96b ; Aharonov:99 ; Zalka:96 ; Gottesman:97b ; Gottesman:99 ; Gottesman:99a ; Preskill:97a ; Steane:99a ; Steane:03 ; Knill:04 ; Reichardt:04 ; Knill:05 :
1. A1*Gates can be executed in a time* $`\tau _g`$ *such that* $`\tau _g\omega =O(\pi )`$, *where* $`\omega `$ *is a typical Bohr or Rabi frequency*.<sup>1</sup><sup>1</sup>1One might object that slower (even adiabatic) gates could be used instead. We analyze this possibility in detail in Section V.4, and show that it does not lead to an improvement.
2. A2*A constant supply of fresh and nearly pure ancillas*: at every time step we are given a supply of many qubits in the state $`|0`$, each of which can be faulty with some error parameter $`\eta 1`$.
3. A3*Error correlations decay exponentially in time and space*.
Some remarks:
(i) A1 is not typically stated explicitly in the FT-QEC literature, but can be understood as resulting from the definition of a quantum gate, which is a unitary transformation $`U=\mathrm{exp}(iA)`$; when $`A=\tau _gH`$, where $`H`$ is a Hamiltonian generating the gate, A1 follows from the absence of a free parameter: when $`\tau _g`$ is scaled up $`H`$ (and hence its eigenvalues) must be scaled down, and vice versa.
(ii) The distinction between Bohr and Rabi frequencies in A1 is related to the application of constant vs periodic control fields, respectively. In the case of a constant control field A1 can be understood as the condition that saturates the โMargolus-Levitin theoremโ Margolus:98 , which states that the time required to transform an initial state $`|\psi `$ to an orthogonal state $`|\psi ^{}`$ using a constant Hamiltonian $`H`$ is lower-bounded by $`\tau _{\mathrm{min}}=\pi \mathrm{}/(2E)`$, where $`E=\psi |H|\psi `$; when $`|\psi `$ is an eigenstate of $`H`$ we have $`\tau _g\pi /(2\omega )`$, where $`\omega =E/\mathrm{}`$ is the Bohr frequency. See also Andrecut:04 for the adiabatic version of the Margolus-Levitin theorem, and Gea:02 for a lower bound on the amount of energy needed to carry out an elementary logical operation on a quantum computer, with a given accuracy and in a given time. In the case of periodic control fields one can understand A1 as the result of the standard solution to the driven two-level atom problem, where the probability of a transition between ground and excited state is given by $`(\mathrm{\Omega }_\mathrm{R}/\mathrm{\Omega }_\mathrm{R}^{})\mathrm{sin}^2(\mathrm{\Omega }_\mathrm{R}^{}t/2)`$, where $`\mathrm{\Omega }_\mathrm{R}`$ is the Rabi frequency and $`\mathrm{\Omega }_\mathrm{R}^{}=(\mathrm{\Omega }_\mathrm{R}+\delta ^2)^{1/2}`$, where $`\delta `$ is the detuning. This expression for the transition probability yields $`\tau _g\mathrm{\Omega }_\mathrm{R}^{}=O(\pi )`$.
(iii) A2 is shown to be necessary in Aharonov:96b . A3 is stated clearly in Aharonov:99 (see the discussion in Sections 2.10 and 10 there). These, and additional assumptions \[such as constant fault rate (independent of number of qubits) and parallelism (to correct errors in all blocks simultaneously)\] are explicitly listed, e.g., also in Dennis:02 , Section II.
(iv) A3 is usually related to the Markovian assumption, however both notions, the space-time correlations of errors and the Markovian property, need some comments and explanations. Using the convolutionless formalism in the theory of open systems (see, e.g., Breuer:02 ) it is always possible to resolve the total superoperator $`\mathrm{\Lambda }(t)`$ as
$$\mathrm{\Lambda }(t)=\underset{i=1}{\overset{n}{}}\mathrm{\Lambda }_iU_i$$
(1)
where $`U_i`$ are ideal unitary superoperators (corresponding to quantum logic gates), and $`\mathrm{\Lambda }_i`$ are linear maps, not necessarily completely positive (CP) or even positive. If $`\mathrm{\Lambda }_i`$ are CP then we can always realize them by coupling to an evironment which is โrenewed each time stepโ. This is the โMarkovian conditionโ as formulated in Aharonov:99 (section 2.10). However, complete positivity is not a necessary condition for QEC, which only requires a linear structure Knill:97b ; ShabaniLidar:06 . To obtain the Threshold Theorem one needs the following bound on the probability Aharonov:99 \[Eq. (2.6)\]:
$$\text{Pr}(\text{fault path with }k\text{ errors})c\eta ^k(1\eta )^{vk},$$
(2)
where $`\eta `$ is the probability of a single error, $`c`$ is a certain constant independent of $`\eta `$, and $`v`$ is the number of error locations in the circuit. This bound implies that the $`k`$-qubit errors should scale as $`\eta ^k`$, i.e., that in the decomposition of $`\mathrm{\Lambda }_j`$ into $`k`$-qubit superoperators $`L_j(k)`$
$$L_j(k)c\eta ^k.$$
(3)
As discussed in Alicki:02 (within the Born approximation), the condition (3) can strictly be satisfied only for temporally *exponentially* decaying reservoir correlation functions, while for realistic reservoir models the temporal decay is generically powerlike. The decay of reservoir correlation functions (i.e., localization in time) translates into localization of errors in space due to the finite speed of error propagation. On the other hand it is widely believed that the Markovian model can be understood as arising, to an excellent approximation, from coupling to a reservoir which is not only renewed at each time step, but whose influence is independent of the actual Hamiltonian dynamics of the open system, and is localized in space (independent errors model) Giulini:book . A large part of the present paper is devoted to a critical discussion of this claim.
(v) We note that the recent papers on FT-QEC theory Terhal:04 ; Aliferis:05 ; Aharonov:05 relax the (Markovian) assumption A3, but do make A1 (implicitly) and A2. We comment on these papers in Section VI.5.
## III Review of Markovian Master Equations
The field of derivations of the quantum Markovian master equation (MME) is strewn with pitfalls: it is in fact non-trivial to derive the MME in a fully consistent manner. There are essentially two types of fully rigorous approaches, known as the *singular coupling limit* (SCL) and the *weak coupling limit* (WCL), both of which we consider below. See, e.g., the books Alicki:87 ; Breuer:book for more details, as well as the derivation in Alicki:89 .
Consider a system and a reservoir (bath), with self Hamiltonians $`H_S^0`$ and $`H_R`$, interacting via the Hamiltonian $`H_{SR}=\lambda SR`$, where $`S`$ ($`R`$) is a Hermitian system (reservoir) operator and $`\lambda `$ is the coupling strength. A more general model of the form $`H_{SR}=_\alpha \lambda _\alpha S_\alpha R_\alpha `$ can of course also be considered and results in the same qualitative conclusions. Thus the total Hamiltonian is
$$H=(H_S^0+H_C(t))I_R+I_SH_R+H_{SR},$$
(4)
where $`H_C(t)`$ describes control over the quantum device (system), and $`I`$ is the identity operator.
The SCL and WCL derivations start from the expansion of the propagator $`\mathrm{\Lambda }`$ for the reduced, system-only dynamics,
$$\rho _S(t)=\mathrm{\Lambda }(t,0)\rho _S(0),$$
(5)
computed in the interaction picture with respect to the renormalized, *physical*, time-dependent Hamiltonian $`H_S(t)=H_S+H_C(t)`$, where
$$H_S=H_S^0+\lambda ^2H_1^{\mathrm{corr}}(t)+\mathrm{}.$$
(6)
The renormalizing terms containing powers of $`\lambda `$ are โLamb-shiftโ corrections due to the interaction with the bath (see, e.g., Lidar:CP01 ). The lowest order (Born) approximation with respect to the coupling constant $`\lambda `$ yields $`H_1^{\mathrm{corr}}`$, while the higher order terms ($`\mathrm{}`$) require going beyond the Born approximation. Introducing a cumulant expansion for the propagator,
$$\mathrm{\Lambda }(t,0)=\mathrm{exp}\underset{n=1}{\overset{\mathrm{}}{}}[\lambda ^nK^{(n)}(t)],$$
(7)
one finds that $`K^{(1)}=0`$. The Born approximation consists of terminating the cumulant expansion at $`n=2`$, whence we denote $`K^{(2)}K`$:
$$\mathrm{\Lambda }(t,0)=\mathrm{exp}[\lambda ^2K(t)+O(\lambda ^3)].$$
(8)
One obtains
$`K(t)\rho _S={\displaystyle _0^t}๐s{\displaystyle _0^t}๐uF(su)S(s)\rho _SS(u)^{}`$
$`+(\mathrm{similar}\mathrm{terms})`$ (9)
as the first term in a cumulant expansion Alicki:89 . Here $`F(s)=\mathrm{Tr}(\rho _RR(s)R)`$ is the autocorrelation function, where $`\rho _R`$ is the reservoir state and $`R(s)`$ is $`R`$ in the $`H_R`$-interaction picture, and $`S(u)`$ is $`S`$ in the interaction picture with respect to the physical Hamiltonian $`H_S(t)`$. The โsimilar termsโ in Eq. (9) are of the form $`\rho _SS(s)S(u)^{}`$ and $`S(s)S(u)^{}\rho _S`$.
At first sight $`K(t)t^2`$, and this is true for small times (Zeno effect Facchi:PRL02 ). The Markov approximation means that we can replace $`K(t)`$ by an expression that is linear in $`t`$, i.e.
$$K(t)_0^t(s)๐s$$
(10)
where $`(t)`$ is a time-dependent Lindblad generator. That the Lindblad generator can be time-dependent even after transforming back to the Schrรถdinger picture is important for our considerations below.
### III.1 Singular Coupling Limit
The SCL approach we present in this subsection underlies the standard derivation of the MME that can be found in almost any text concerning the Markov approximation, though not always under the heading โSCLโ (e.g., Carmichael:book , p.8, Eq. (1.36)). The rigorous derivation of the SCL is briefly discussed (with references) in Alicki:87 , pp.36-38. It is based on a rescaling of the bath and system-bath Hamiltonians, which physically makes sense in the high-temperature limit only. We will shortly see the emergence of this limit.
In essence, the โnaive SCL-Markov approximationโ is obtained by the ansatz $`F(s)=a\delta (s)`$ for the autocorrelation function, whence
$$L(s)\rho _S=aS(s)\rho _SS(s)^{}+(\mathrm{similar}\mathrm{terms}).$$
(11)
As a consequence, return to the Schrรถdinger picture gives a MME with the dissipative part independent of the Hamiltonian:
$`{\displaystyle \frac{d\rho _S}{dt}}`$ $`=`$ $`i[H_S(t),\rho _S]+\rho _S,`$
$`\rho _S`$ $``$ $`{\displaystyle \frac{1}{2}}\lambda ^2a[S,[S,\rho _S]]`$ (12)
More precisely, we must consider the multi-time bath correlation functions $`F(t_1,\mathrm{},t_n):=\mathrm{Tr}[\rho _RR(t_1)\mathrm{}R(t_n)]:=R(t_1)\mathrm{}R(t_n)`$. Here $`R(t):=\mathrm{exp}(iH_Rt)R\mathrm{exp}(iH_Rt)`$ are the bath operators in the interaction picture, $`\rho _R=\mathrm{exp}(\beta H_R)/Z`$ (where $`\beta =1/kT`$, $`Z=\mathrm{Tr}[\mathrm{exp}(H_R/kT)]`$) is the bath thermal equilibrium state at temperature $`T`$, which is a stationary state of the reservoir, i.e., $`[H_R,\rho _R]=0`$. The influence of the environment on the system is entirely encoded into the $`\{F(t_1,\mathrm{},t_n)\}_{n=2}^{\mathrm{}}`$.<sup>2</sup><sup>2</sup>2$`F(t_1)`$ is constant by stationarity. We reserve the notation $`F(t)`$ for $`F(t_1,t_2)F(t_2t_1)`$ below. Heuristically, the Markov approximation can be justified under the following conditions:
1. The lowest order correlation function,
$$F(t)=R(s+t)R(s)=_{\mathrm{}}^{\mathrm{}}G(\omega )e^{i\omega t}๐\omega ,$$
(13)
can be approximated by a Dirac delta function:<sup>3</sup><sup>3</sup>3Note that stationarity implies that $`F(t)`$ does not depend on $`s`$.
$$F(t)\left(_{\mathrm{}}^{\mathrm{}}F(s)๐s\right)\delta (t)=G(0)\delta (t)$$
(14)
(white-noise approximation). Eq. (13) defines the *spectral density* $`G(\omega )`$, which is a key object in the theory.
2. Higher order correlation functions exhibit a Gaussian-type behavior, i.e., can be estimated by sums of products of the lowest order ones, and then, by condition (14), decay sufficiently rapidly.
Let us now comment on the physical relevance of the white-noise approximation.
First, the condition (14) cannot be satisfied in general. For example, in the important case of linear coupling to a bosonic field (e.g., electromagnetic field, phonons in solid state), we have $`G(0)=0`$, which means (by inverse Fourier transform) that $`_{\mathrm{}}^+\mathrm{}F(t)๐t=0`$, and therefore *$`F(t)`$ cannot be well approximated by* $`\delta (t)`$.
Second, even for models with $`G(0)>0`$ there exists a universal relation, the so-called Kubo-Martin-Schwinger (KMS) condition, $`R(t)R(0)=R(0)R(t+i\beta )`$, which is valid for all quantum systems at thermal equilibrium. This implies:
$$G(\omega )=e^{\beta \omega }G(\omega ).$$
(15)
(See, e.g., Alicki:87 \[pp.90-91\], Thirring:book \[pp.176-177\], or Breuer:book \[p.137\].) The fundamental importance of the KMS condition is captured by the fact that it is necessary in order for thermodynamics to hold. The KMS condition implies a strong asymmetry of the spectral density $`G(\omega )`$ for low $`T`$, where $`T`$ is measured relative to the presence of $`kT`$ energy scales in the bath, i.e., relative to the range where $`G(\omega )`$ is non-vanishing. The KMS condition is relevant to our discussion since we make the reasonably minimalistic assumption that the reservoir (not the QC) is in thermal equilibrium.<sup>4</sup><sup>4</sup>4One may challenge the notion that the bath must always be in thermal equilibrium. E.g., consider an atom in a microwave cavity, with the cavity electromagnetic field initially in thermal equilibrium. Now suppose the atom is driven and is coupled to the cavity electromagnetic field, which therefore is no longer in equilibrium. However, is the internal electromagnetic field the relevant environment, or is it the external one? Clearly, the electromagnetic field inside the cavity is not a reservoir but itself a part of the system. This is because: a) its spectrum is discrete, b) its coupling to the atom (close to resonance) is enhanced. The reason these considerations matter is because b) implies the strong coupling regime, hence failure of the initial state tensor product structure assumption, hence difficulties with the separation of the system from the reservoir (dressed atom picture); a) implies that $`F(t)`$ is (quasi)-periodic with short Poincarรฉ recurrences, hence a strongly non-Markovian regime, and thus associated difficulties for Markovian FT-QEC. On the other hand the external electromagnetic field has a continuous spectrum and the state product structure is easily satisfied, hence qualifies as a reservoir. This example merely serves to illustrate accepted notions regarding the division into well defined system and bath; for most practical purposes a thermal equlibrium is the simplest and most relevant model of an environment, and FT-QEC theory must be applicable to this setting.
Third, $`G(\omega )`$ need not be flat even at high $`T`$ (indeed, the KMS condition only implies that $`G(\omega )`$ is symmetric at high $`T`$). For example, this is the case for the electromagnetic field and for phonons, for which at $`T>0`$ one has $`G(\omega )\omega ^3/(1e^{\mathrm{}\omega /kT})`$ for $`|\omega |\omega _{\mathrm{cut}}`$, and $`G(\omega )=0`$ for $`|\omega |>\omega _{\mathrm{cut}}`$. One can see that for high $`T`$ ($`kT\mathrm{}\omega _{\mathrm{cut}}`$), $`G(\omega )kT\omega ^2`$ is symmetric. Here $`\omega _{\mathrm{cut}}`$ is the Debye frequency in the case of phonons, while for the electromagnetic field $`\omega _{\mathrm{cut}}`$ should tend to infinity in the renormalization procedure. A flat $`G(\omega )`$ means a structureless bath, while physical systems always have a nontrivial structure depending on relevant energy scales.<sup>5</sup><sup>5</sup>5It is interesting to note that even if we try to enforce a flat $`G(\omega )`$ by, e.g., choosing an appropriate form factor for the spin-boson system, the obtained model โ the so-called โOhmic caseโ โ is mathematically and physically ill defined (see Alicki:02a ).
Now let us return to the implications of the SCL assumptions for the problem of FT-QEC. In order to derive the SCL from first principles, one rescales $`H_RH_R/ฯต^2`$, rescales $`H_{SR}H_{SR}/ฯต`$, but keeps $`H_S`$ and $`\rho _R`$ fixed.<sup>6</sup><sup>6</sup>6Note that because different Hamiltonians are rescaled differently, this rescaling procedure is *not* equivalent to a direct rescaling of the time variable (which is what is done in the WCL, below). The idea of this rescaling is that it accelerates the reservoirโs evolution (via $`H_RH_R/ฯต^2`$) and hence produces faster decay of the reservoir correlations, $`F(t)`$. To see this, note that the rescaling $`H_{SR}H_{SR}/ฯต`$ increases the amplitude $`F(0)`$ to $`F(0)/ฯต^2`$ (proportional to $`H_{SR}^2`$), while keeping the strength of the noise $`_{\mathrm{}}^+\mathrm{}F(t)๐t=G(0)`$ fixed (as can be seen via a change of variables $`tt/ฯต^2`$ in the integral). This implies a faster decay of $`F(t)`$. The rescaling procedure is specifically designed to yield the delta correlation \[Eq. (14)\] in the limit as $`ฯต0`$. Note that if $`\rho _R`$ is at thermal equilibrium at temperature $`T`$ with respect to $`H_R`$, then, since $`\rho _R=\mathrm{exp}(\beta H_R)/Z`$ is fixed, it must be at thermal equilibrium with respect to $`H_R/ฯต^2`$ at the temperature $`T/ฯต^2\mathrm{}`$, whence our mention of the high temperature limit, above. Further note that $`H_S`$ is not rescaled since the SCL is (artificially) designed to produce โwhite noiseโ on the natural time scale of systemโs evolution, which is given by $`H_S`$.
Another, equivalent way to understand the emergence of the high-$`T`$ limit is the following: For the Markovian condition $`F(t)a\delta (t)`$ to hold the spectral density must be flat: $`G(\omega )=\mathrm{const}`$. However, this is possible only in the limit $`T\mathrm{}`$ of the KMS condition. More precisely, *the Markovian condition can hold only if* $`kT\omega `$ *over the entire spectrum of the systemโs Bohr frequencies*. Strictly speaking, $`G(\omega )`$ is never constant. The variation of $`G(\omega )`$ happens over the โthermal memoryโ time $`\tau _{\mathrm{th}}:=1/kT`$. In the infinite $`T`$ limit we then recover the case of zero memory-time, i.e., Markovian dynamics. Physically, it is enough to assume that $`G(\omega )`$ is essentially constant over the interval $`[\omega _0,\omega _0]`$ where $`kT>\omega _0`$ systemโs Bohr frequencies. I.e., system energy scales must be compared to $`1/\tau _{\mathrm{th}}`$ and this leads to the important realization that *the Markovian approximation can be consistent with the KMS condition only in the high temperature regime* $`kTE`$, *where* $`E`$ *is the system energy scale*. As we argue below, this fact presents a serious difficulty in the context of FT-QEC, the issue being essentially that the requirement of a constant supply of nearly pure and cold ancillas contradicts the high-$`T`$ limit needed for the Markov approximation to hold.
### III.2 Weak Coupling Limit
In the SCL approach above there was no restriction on the time-dependence of the system Hamiltonian. However, the price paid is the high-$`T`$ limit. Moreover, while mathematically the SCL is rigorous in the scaling limit, it is inconsistent with thermodynamics except in the $`T\mathrm{}`$ limit. On the other hand, the derivation by Davies, in his seminal 1974 paper Davies:74 , is perhaps the only derivation of the MME that is entirely consistent from *both the mathematical and physical* points of view. The Davies approach is based not on the high-$`T`$ limit, but rather on the physically plausible idea of weak coupling. This is natural and consistent with thermodynamics at all temperatures.
More specifically, Daviesโ derivation does not invoke a flatness condition on $`G(\omega )`$ but is, of course, still subject to the KMS condition. In the Davies approach the Markov approximation is a consequence of weak coupling (and hence slow dynamics of the system in the interaction picture), and time coarse-graining, which leads to cancellation of the non-Markovian oscillating terms. The price we pay is the invalidity of this approach for time-dependent Hamiltonians, except in the adiabatic case. We explain this important comment below. Hence, while the Davies approach does not require the high-$`T`$ limit, it imposes severe restrictions on the speed of quantum gates.
In his rigorous derivation Davies replaced the heuristic condition (14) by the weaker
$$|F(t)|๐t<\mathrm{}.$$
(16)
This condition avoids the difficulties originating from the singularity of the SCL condition (14), and preempts the corresponding problems with the high-$`T`$ limit.<sup>7</sup><sup>7</sup>7In some sense the weak coupling limit is similar to the Central Limit Theorem (CLT) in probability, and condition (16) is analogous to a rough upper bound on the second moment in the CLT. If it is not satisfied then the noise may be not Gaussian in the weak coupling limit. The value of $`|F(t)|๐t`$ itself does not provide any meaningful physical parameter and can depend on some regularization/cut-off parameters.<sup>,</sup><sup>8</sup><sup>8</sup>8One can go further and ask how generic the Markovian case is, in the sense that Eq. (16) is satisfied. In fact, typically $`F(t)`$ decays as $`1/t^\alpha `$ (e.g., for the vacuum bath $`\alpha =4`$ Alicki:02 ), which means that in some cases ($`\alpha 1`$) Eq. (16) can be violated. For a systematic treatment of these non-Markovian effects see, e.g., Breuer:02 . We now consider the cases of a constant, periodic, and arbitrarily time-dependent control Hamiltonian. The constant case is the one originally treated by Davies Davies:74 , and extended in Davies:78 to time-dependent Hamiltonians assuming a slow (โadiabaticโ) change on the dissipation time scale $`\lambda ^2t`$. The non-constant cases we study here have, as far as we know, not been published before in the general scientific literature.
#### III.2.1 WCL for Constant $`H_C`$: Summary of the Original Davies Derivation
We present a simplified version of the discussion of the Markov approximation in Alicki:89 . Denote by $`E_k`$ the Bohr energies (eigenvalues of $`H_S`$), let $`\omega \{\omega _{kl}=E_kE_l\}_{k,l}`$, and let $`S_\omega `$ be the discrete Fourier components of the interaction picture $`S`$, i.e.,
$$S(t)=\mathrm{exp}(iH_St)S\mathrm{exp}(iH_St)=\underset{\omega }{}S_\omega \mathrm{exp}(i\omega t),$$
(17)
where $`H_S`$ is the renormalized (physical) system Hamiltonian: the sum of the โbareโ $`H_S^0`$ and a Lamb shift term (bath induced), as in Eq. (6). Equivalently,
$$[H_S,S_\omega ]=\omega S_\omega .$$
(18)
We remark that in the original Davies paper the Bohr energies and Eq. (18) are computed with respect to the bare Hamiltonian $`H_S^0`$. Here we use the physical Hamiltonian $`H_S`$ in order to take into account the fact that the Lamb shift term, although formally proportional to $`\lambda ^2`$, can be large or even infinite after cut-off removal.
Then, it follows from Eq. (9) that
$`K(t)\rho _S`$ $`=`$ $`{\displaystyle \underset{\omega ,\omega ^{}}{}}S_\omega \rho _SS_\omega ^{}^{}{\displaystyle _0^t}e^{i(\omega \omega ^{})u}๐u{\displaystyle _u^{tu}}F(\tau )e^{i\omega \tau }๐\tau `$ (19)
$`+`$ $`(\mathrm{similar}\text{ }\mathrm{terms}).`$
The weak coupling limit is next formally introduced by rescaling the time $`t`$ to $`t/\lambda ^2`$ (van Hove limit). This enables two crucial approximations, which are valid in the resulting large-$`t`$ limit:
1. We replace<sup>9</sup><sup>9</sup>9In a more rigorous treatment the Cauchy principal value must be used, but the result is essentially the same Alicki:89 .
$$_0^te^{i(\omega \omega ^{})u}๐ut\delta _{\omega \omega ^{}}.$$
(20)
This makes sense for
$$t\mathrm{max}\{1/(\omega \omega ^{})\}.$$
(21)
*This violates A1, expressed in terms of the Bohr frequencies*. We see here already the emergence of an adiabatic criterion for the validity of the Markov approximation.
2. We replace $`_u^{tu}F(\tau )e^{i\omega \tau }๐\tau `$ by the Fourier transform:
$$_u^{tu}F(\tau )e^{i\omega \tau }๐\tau G(\omega )=_{\mathrm{}}^{\mathrm{}}F(\tau )e^{i\omega \tau }๐\tau .$$
(22)
The physical validity of the last approximation is usually ignored, though one can make the following argument: On the LHS of Eq. (22), for a given Bohr frequency $`\omega `$ the Fourier-like integral must sample the function $`F(\tau )`$ with sufficiently high accuracy so that the Fourier transform approximation will be valid. To this end one needs a time $`t`$ such that: (i) $`t1/\omega `$. This is a weaker condition than the previous one \[$`t\mathrm{max}\{1/(\omega \omega ^{})\}`$\] which involves differences of Bohr frequencies. (ii) The time $`t`$ must be also much longer than the time scale of the wildest variations of $`F(\tau )`$, which is typically \[as may be checked for simple models of spectral densities $`G(\omega )`$\] given by $`1/\omega _{\mathrm{cut}}`$, where $`\omega _{\mathrm{cut}}`$ is a high-frequency cutoff. When $`\omega <\omega _{\mathrm{cut}}`$ (i) implies (ii). Therefore typically Eq. (20) is a stronger assumption than Eq. (22).
Applying the approximations (20) and (22), we obtain $`K(t)\rho _S=t_\omega S_\omega \rho _SS_\omega ^{}G(\omega )+(\mathrm{similar}`$ $`\mathrm{terms})`$, and hence it follows from Eq. (10) that $`(s)=`$ is the Davies generator in the familiar Lindblad form:
$`{\displaystyle \frac{d\rho _S}{dt}}`$ $`=`$ $`i[H_S,\rho _S]+\rho _S,`$
$`\rho _S`$ $``$ $`{\displaystyle \frac{1}{2}}\lambda ^2{\displaystyle \underset{\omega }{}}G(\omega )([S_\omega ,\rho S_\omega ^{}]+[S_\omega \rho ,S_\omega ^{}])`$ (23)
Several remarks are in order:
(i) The absence of off-diagonal terms in Eq. (23), compared to Eq. (19), is the hallmark of the Markovian limit. Namely, the Davies derivation relies on the cancellation of the non-Markovian off-diagonal terms $`_{\omega \omega ^{}}S_\omega \rho _SS_\omega ^{}^{}_0^te^{i(\omega \omega ^{})u}๐u`$. This time coarse-graining is possible due to integration over the fast-oscillating $`_0^te^{i(\omega \omega ^{})u}`$ terms over a long timescale, i.e., over $`t\mathrm{max}\{1/(\omega \omega ^{})\}`$ (see also Lidar:CP01 ). As remarked above, this violates A1, expressed in terms of the Bohr frequencies.
(ii) It follows from Bochnerโs theorem applied to the Fourier transform definition of $`G(\omega )`$ that $`G(\omega )0`$ Alicki:87 \[p.90\], Breuer:book \[p.136\]; this result is essential for the complete positivity of the Markovian master equation in the WCL.
(iii) Daviesโ derivation showed implicitly that the notion of โbathโs correlation timeโ is not well-defined โ Markovian behavior involves a rather complicated cooperation between system and bath dynamics. More specifically, the relations (23) and (18) together imply that the noise and $`H_S`$ are strongly correlated. In other words, contrary to what is often done in phenomenological treatments, *one cannot combine arbitrary* $`H_S`$*โs with given* $`S_\omega `$*โs.* This point is particularly relevant in the context of FT-QEC, where it is common to assume Markovian dynamics and apply arbitrary control Hamiltonians.
Davies did not consider time-dependent system Hamiltonians in Davies:74 , but it is possible to generalize his derivation to allow for slowly varying system Hamiltonians Davies:78 ; Alicki:79 ; Alicki:89 . That is, whenever the time scale of the variation of $`H_C(t)`$ is much longer than the inverse of the typical Bohr frequency (of $`H_S`$), it is possible to add $`H_C(t)`$ to the system Hamiltonian in Eq. (23), necessitating at the same time this change also in Eq. (18). This is a type of adiabatic limit (indeed, the $`S_\omega `$ in Eq. (18) can be interpreted, with $`H_S`$ replaced by $`H_S+H_C(t)`$, as being adiabatic eigenvectors of the superoperator $`[H_S+H_C(t),]`$). We note that an alternative approach to adiabaticity in open quantum systems was recently developed in Ref. SarandyLidar:04 . This approach, while being very general, is more phenomenological in that it postulates a convolutionless master equation, and then derives corresponding adiabaticity conditions. Closer in spirit to the Davies derivation is another recent approach to adiabaticity in open systems, which assumes slow system variation together with weak system-bath coupling Thunstrom:05 .
#### III.2.2 WCL for Periodic Driving: Floquet Analysis
Before considering the case of periodic $`H_C`$ let us consider briefly once more the case of a constant Hamiltonian in the so-called covariant dissipation setting. Covariance is the commutation condition $`=`$ where $`=[H_S,]`$ is the super-operator constant Hamiltonian, and $``$ is the Davies generator. Covariance is an abstract property which is automatically fulfilled for the Davies generator.<sup>10</sup><sup>10</sup>10This can be verified by directly computing $``$ and making use of Eq. (18) and the relation $`[A,BC]=[A,B]C+B[A,C]`$ (for operators $`A,B`$ and $`C`$). A more elegant way to see this is to consider $`(t)=\mathrm{exp}(it)\mathrm{exp}(it)`$ and note that Eq. (18) implies that $`S(t)`$ and $`S^{}(t)`$ rotate in opposite directions. Hence $`(t)=`$, whence $`d(t)/dt=0`$ gives the result. It is convenient since it implies factorization of the full propagator into Hamiltonian and dissipative parts. Markovian dynamics obtained in the WCL as discussed above takes the form
$$\frac{d\rho }{dt}=(i+)\rho ,t0,$$
(24)
where the most general form of the Lindblad (or Davies) $``$ satisfying Eq. (24) is
$$\rho =\frac{1}{2}\underset{\{\omega \},j}{}\left([V_j(\omega ),\rho V_j(\omega )^{}]+[V_j(\omega )\rho ,V_j(\omega )^{}]\right).$$
(25)
Here $`\{\omega \}\mathrm{Spectrum}()`$, i.e., the Bohr frequencies (differences of eigenvalues of $`H`$), and
$$V_j(\omega )=\omega V_j(\omega )$$
(26)
\[i.e., Eq. (18)\]. The solution, i.e., the dynamical semigroup is
$$\rho (t)=\mathrm{\Lambda }(t)\rho (0),\mathrm{\Lambda }(t)=e^{it}e^t.$$
(27)
Now consider a periodic control Hamiltonian with period $`\mathrm{\Theta }`$
$$H_C(t)=H_C(t+\mathrm{\Theta }),\mathrm{\Omega }=2\pi /\mathrm{\Theta }.$$
(28)
(Note that $`\mathrm{\Omega }`$ is *not* the Rabi frequency, which throughout this paper we denote by $`\mathrm{\Omega }_\mathrm{R}`$.) The situation is then very similar to the standard (time-independent $`H_C`$) WCL, but the set of โeffective Bohr frequenciesโ (Floquet spectrum) $`\omega `$ is now larger and is of the form $`\{\omega +q\mathrm{\Omega }\}`$, $`q=0,\pm 1,\mathrm{}`$. Here $`\omega `$ are Bohr frequencies for the Floquet unitary \[defined in Eq. (30) below\], i.e., differences of eigenvalues $`ฯต_\alpha `$ of the Floquet unitary, rather than $`\{\omega \}=\mathrm{Spectrum}()`$ as above. As this set of โeffective Bohr frequenciesโ is discrete the WCL still works, but the final Davies generator is more complicated, as we now show.
Define the time-ordered unitary propagator
$$U(t,s)๐ฏ\mathrm{exp}\left(i_s^tH_S(u)๐u\right),ts$$
(29)
which satisfies the properties $`U(s,t)U(t,s)^1=U(t,s)^{}`$, $`U(t,s)U(s,u)=U(t,u)`$, $`U(t+\mathrm{\Theta },s+\mathrm{\Theta })=U(t,s)`$, and $`\frac{d}{dt}U(t,s)=iH_S(t)U(t,s)`$, $`\frac{d}{dt}U(t,s)^{}=iU(t,s)^{}H_S(t)`$. The *Floquet unitary operator* is
$$F(s)U(s+\mathrm{\Theta },s)\text{},$$
(30)
with corresponding super-operator action
$$(s)\rho F(s)\rho F(s)^{},$$
(31)
and Floquet eigenvectors $`|\varphi _\alpha `$ and eigenvalues (quasi-energies) $`ฯต_\alpha `$ satisfying<sup>11</sup><sup>11</sup>11Note that the Floquet Hamiltonian $`H_S(t)id/dt`$ operates on a different Hilbert space than $`F(0)`$ (the space of periodic functions with values in the systemโs Hilbert space). But its eigenvalues coincide with $`ฯต_\alpha `$ from Eq. (32).
$$F(0)|\varphi _\alpha =e^{iฯต_\alpha \mathrm{\Theta }}|\varphi _\alpha .$$
(32)
It follows from standard Floquet theory that
$$U(t,0)|\varphi _\alpha =e^{itฯต_\alpha }\underset{q๐}{}e^{itq\mathrm{\Omega }}|\varphi _\alpha (q),$$
(33)
i.e., the set $`\{|\varphi _\alpha (q)\}`$ is a complete basis. Therefore we have at most as many $`q`$โs as the dimension of the Hilbert space. That the number of $`q`$โs is finite is important for our considerations below.
We call a Lindblad generator $``$ a โcovariant dissipative perturbation of $`H_S(t)`$โ if
$$(0)=(0)$$
(34)
We will assume this property, similarly to the case of a constant Hamiltonian described above. In fact, covariance holds for a periodic $`H_S(t)`$ and also for the corresponding WCL Davies generator. One can then derive the *covariant master equation* (we sketch the derivation below):
$$\frac{d\rho }{dt}=\left(i(t)+(t)\right)\rho ,t0,$$
(35)
\[compare to Eq. (24)\] where
$`(t)`$ $`=`$ $`๐ฐ(t,0)๐ฐ(t,0)^{},`$ (36)
$`{\displaystyle \frac{d}{dt}}๐ฐ(t,s)`$ $`=`$ $`i(t)๐ฐ(t,s),`$ (37)
and where the general form of $``$ appearing in Eq. (36) is given by Eq. (25), with $`V_j(\omega )`$ now being eigenvectors of $`(0)`$,
$$(0)V_j(\omega )=e^{i\omega \mathrm{\Theta }}V_j(\omega ),$$
(38)
rather than of $``$, as in Eq. (26). Moreover, here $`\{\omega \}\{ฯต_\alpha ฯต_\beta \}`$, where $`ฯต_\alpha `$ are quasi-energies (effective Bohr frequencies) of the Floquet operator, rather $`\{\omega \}\mathrm{Spectrum}()`$ as we saw in the case of constant $`H_C`$.
The solution replacing Eq. (27) is
$`\rho (t)`$ $`=`$ $`\mathrm{\Lambda }(t,s)\rho (s),\text{ }ts`$
$`\mathrm{\Lambda }(t,s)`$ $`=`$ $`๐ฏ\mathrm{exp}\left\{{\displaystyle _s^t}\left(i(u)+(u)\right)๐u\right\}`$ (39)
By direct computation one can prove the following properties:
$`(t+\mathrm{\Theta })`$ $`=`$ $`(t),`$ (40)
$`(s)(s)(s)^{}`$ $`=`$ $`(s),`$ (41)
$`\mathrm{\Lambda }(t,s)\mathrm{\Lambda }(s,u)`$ $`=`$ $`\mathrm{\Lambda }(t,u)\text{ for }tsu,`$ (42)
$`\mathrm{\Lambda }(t+\mathrm{\Theta },s+\mathrm{\Theta })`$ $`=`$ $`\mathrm{\Lambda }(t,s),`$ (43)
$`\mathrm{\Lambda }(t,s)`$ $`=`$ $`๐ฐ(t,s)e^{(ts)(s)}.`$ (44)
To derive the covariant master equation (35) one considers the standard picture of an open system $`S+R`$ with the total Hamiltonian
$$H_{SR}(t)=H_S^0(t)+H_R+\underset{k}{}S_kR_k,$$
(45)
(we neglect the Lamb shift correction here; it can be included, changing $`H_S^0(t)`$ into the physical Hamiltonian $`H_S(t)`$, by a suitable renormalization procedure), stationary reservoir state $`\rho _R`$, $`[H_R,\rho _R]=0`$, $`\mathrm{Tr}(\rho _R)_R`$, $`R_k_R=0`$. Then, exactly following a Davies-like calculation using a Fourier decomposition of $`S(t)`$, now governed by a periodic Hamiltonian, and making in particular again the crucial assumption Eq. (21), which now reads
$$t\mathrm{max}\{1/(\omega \omega ^{}+m\mathrm{\Omega })\},m=0,\pm 1,\pm 2,\mathrm{}$$
(46)
with $`|m|`$ upper-bounded by the dimension of the Hilbert space \[see remark after Eq. (33)\], one obtains Eq. (35) in the Davies WCL. The explicit form of the generator is:
$`\rho `$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{k,l}{}}{\displaystyle \underset{q๐}{}}{\displaystyle \underset{\{\omega \}}{}}\widehat{R}_{kl}(\omega +q\mathrm{\Omega })\{[S_l(q,\omega )\rho ,S_k(q,\omega )^{}]`$ (47)
$`+`$ $`[S_l(q,\omega ),\rho S_k(q,\omega )^{}]\}.`$
Here $`\{\omega \}\{ฯต_\alpha ฯต_\beta \}`$, the Floquet spectrum, and
$$\widehat{R}_{kl}(x)=_{\mathrm{}}^{\mathrm{}}e^{itx}R_k(t)R_l_R๐t$$
(48)
$$S_k(q,\omega )=\underset{p๐}{}\underset{\{ฯต_\alpha ฯต_\alpha ^{}=\omega \}}{}\varphi _\alpha (p+q)|S_k|\varphi _\alpha ^{}(p)|\varphi _\alpha \varphi _\alpha ^{}|.$$
(49)
$`S_k(q,\omega )`$ is the part of $`S(t)`$ which rotates with frequency $`\omega +q\mathrm{\Omega }`$ and can be computed using Eq. (33). Note that by diagonalizing the matrices $`\widehat{R}_{kl}`$ one can transform the generator $``$ of Eq. (47) into the form of Eq. (25), which allows one to read off the operators $`V_j(\omega )`$appearing there.
Now to some important comments:
*Timescale analysis*: Note that for the periodic case the differences of โBohr frequenciesโ may be of the order of $`1/\mathrm{\Theta }`$. Hence we conclude from Eq. (46) that one must average over many periods $`\mathrm{\Theta }`$, i.e., require $`t\mathrm{\Theta }`$. This can be interpreted as a condition that โthe environment must learn that the Hamiltonian is periodicโ. This is exactly analogous to the adiabaticity condition in the adiabatic case: $`H(t)`$ must be constant over many inverse Bohr frequencies to โbe recognisedโ by the environment. The periodic WCL is also a coarse-grained time description with the additional time scale $`\mathrm{\Theta }`$. *Note that arbitrarily fast periodic driving (small $`\mathrm{\Theta }`$) is incompatible even with the kind of generalized, finitely localized MME derived here, since then differences of Bohr frequencies matter in Eq. (46)* (recall that $`\mathrm{max}|m|`$ is bounded by the โ typically small โ dimension of the system Hilbert space).
*Where is the Rabi frequency?* Note the dependence of the operators $`S_k(q,\omega )`$ on the Floquet eigenvalue differences $`ฯต_\alpha ฯต_\alpha ^{}`$. The usual Rabi frequency, $`\mathrm{\Omega }_\mathrm{R}=2dE/\mathrm{}`$ ($`d`$ is the dipole moment, $`E`$ is the electric field amplitude) arises in the dipole approximation, which we have not made here. The usual Rabi frequency is replaced in our non-perturbative treatment (in the sense of no multipole expansion) by the difference of Floquet eigenvalues $`ฯต_\alpha ฯต_\alpha ^{}`$ in Eq. (49).<sup>12</sup><sup>12</sup>12One can see that such a term also arises in the usual dipole approximation by considering, e.g., Eq. (2.94) in Carmichael:book . The interaction picture raising and lowering operators $`\sigma _\pm (t)`$ (for a two-level atom driven by a classical field) there oscillate with three โBohr frequenciesโ $`\omega _A,\omega _A\pm \mathrm{\Omega }`$ , where $`\mathrm{\Omega }=2dE/\mathrm{}`$ denotes the usual Rabi frequency. Hence the Rabi frequency is a difference of two Bohr frequencies.
*More on the Rabi frequency*: As we saw, the non-Markovian terms vanish because of the time coarse-grained description. To attain this, we must average over times $`t`$ $`\mathrm{max}_{\omega ,\omega ^{}}\{1/(\omega \omega ^{})\}`$, but must also keep in mind that the longest relevant scale for coarse-graining is given by the exponential decay time $`\tau `$ (a *derived* quantity), i.e., we must have $`t<\tau `$. The Rabi frequency $`\mathrm{\Omega }_\mathrm{R}`$ is a difference of two Bohr frequencies $`\omega ,\omega ^{}`$$`\mathrm{\Omega }_\mathrm{R}=\omega \omega ^{}`$. This implies that coarse-graining does not makes sense if $`\mathrm{\Omega }_\mathrm{R}\tau 1`$ \[since then $`t<\tau 1/\mathrm{\Omega }_\mathrm{R}=1/(\omega \omega ^{})`$, in contradiction to the fundamental requirement on $`t`$\]. In physical terms this means that the width of the spectral line ($`\gamma =1/\tau `$) is larger than the level splitting $`\mathrm{\Omega }_\mathrm{R}`$ (see, e.g., Fig. 2.5 (i),(ii) in Carmichael:book for an illustration in the case of the incoherent fluorescence spectrum) and therefore โthe environment has no time to recognize the details of the spectrumโ. On the other hand, when $`\mathrm{\Omega }_\mathrm{R}\tau 1`$ (not inconsistent with the WCL), $`\mathrm{\Omega }_\mathrm{R}`$ must appear in the generator, as appears from our treatment of the case of periodic driving in the WCL, above. Unfortunately there are examples in the literature where an MME is written down subject to $`\mathrm{\Omega }_\mathrm{R}\tau 1`$ but $`\mathrm{\Omega }_\mathrm{R}`$ does not appear in the generator \[e.g., Eq. (2.96) in Carmichael:book , where $`\mathrm{\Omega }_\mathrm{R}10^{10}\mathrm{Hz}`$ and $`\tau 10^8\mathrm{s}`$\].
*Quantum optics considerations*: The Markov approximation is commonly accepted as an excellent approximation in quantum optics; see, e.g., the discussion of resonance fluorescence in Carmichael:book \[Ch.2\]. This is also the basis for substantial confidence in the possibility of FT-QEC in quantum optical systems, such as trapped ions Cirac:95 and atoms trapped in microwave cavities Turchette:95 . Such arguments are based on the relative flatness of the damping constants $`\gamma (\omega )`$ as a function of frequency. This argument is closely related to the notion of the flatness of the spectral density $`G(\omega )`$ in the SCL, since the damping constants are proportional to $`G(\omega )`$ \[see Eq. (23)\]. For example, below Eq. (2.95) in Carmichael:book the author argues that one can write down a Rabi frequency-independent MME for resonance fluorescence since $`\gamma (\omega _A)`$ and $`\gamma (\omega _A\pm \mathrm{\Omega }_\mathrm{R})`$ (where $`\omega _A`$ is the Bohr frequency) differ by less than 0.01$`\%`$ at optical frequencies and reasonable laser intensities. However, this ignores the corrections due to the Rabi frequency to the operators $`S_k(q,\omega )`$ \[Eq. (49)\]. This disagreement can be traced to the question of at which point in the derivation it is safe to neglect $`\mathrm{\Omega }_\mathrm{R}`$; in Carmichael:book this is done on the basis of the flatness of $`\gamma (\omega )`$ before โa lot of tedious algebraโ Carmichael:book \[p.48\], but our Floquet analysis shows that, in fact, one cannot neglect the Rabi frequency relative to the Bohr frequency. This is relevant for our general discussion since the โ*finitely localized* MMEโ which is the outcome of the Floquet analysis (see next comment) actually exhibits a weak non-Markovian character. Such deviations are, of course, important for FT-QEC, even if the effects are small. We revisit this point in Section V below.
*Are there any non-Markovian effects at work here?* It seems that one should accept the *generalized notion* of a quantum Markovian master equation as the one given by Eqs. (35), (36) and (25), i.e., a master equation with a possibly time-dependent Lindblad generator. In Daviesโ generalization to the time-dependent case Davies:78 (โadiabatic WCLโ) the dissipative generator $`(t)`$ depends on the Hamiltonian at the *same time* $`t`$. This is a type of โ*local* generalized MMEโ. On the other hand, in the periodic WCL treated here, the dissipative generator $`(t)`$ depends on the Hamiltonians $`H_S(u)`$ from an interval, say $`[0,t]`$ ($`t<\mathrm{\Theta }`$), as can be seen from Eq. (36), which involves $`๐ฐ(t,0)`$. This is therefore a type of โ*finitely localized* MMEโ, though one could argue that it exhibits a weakly non-Markovian character because of this dependence of the dissipative generator on the past. On the other hand, a non-Markovian master equation (in the convolutionless formalism Breuer:book ) is also given by Eq. (35), but the generator is *not* of Lindblad form \[in particular, it is not of the form (32)\], and may depend on the Hamiltonian in the *distant* past. The weight of distant past contributions depends on the decay properties of $`F(t)`$ which are, generically, not exponential but rather powerlike. In the WCL the non-Lindbladian terms vanish due to the oscillating character of the $`e^{i(\omega \omega ^{})u}`$ terms in Eq. (19).
*The original Davies derivation*: We note that the Davies result is a limit theorem which states that for a sufficiently small coupling constant the WCL semigroup is a good approximation to the real dynamics. However, Daviesโ theorem itself does not provide the conditions under which a given physical coupling is โsmall enoughโ. In particular, one cannot extract from Daviesโ theorem under what conditions the fast oscillating terms vanish. This can, however, be done by a more heuristic analysis, as done above.
#### III.2.3 WCL for an Arbitrary Pulse
We now consider the case
$$H_C(t)=H_0+f(t)H_1,$$
(50)
i.e., an arbitrary driving field. This is, of course, the case of most interest in FT-QEC. It follows from Fourier analysis that this case can be treated qualitatively as a โsuperpositionโ of periodic perturbations discussed above. For a single frequency $`\mathrm{\Omega }`$ the validity of the Markovian approximation is restricted by the condition (46): $`t\mathrm{max}\{1/(\omega \omega ^{}+m\mathrm{\Omega })\}`$. The discreteness of the frequencies $`\{\omega \}`$ and $`\{m\mathrm{\Omega }\}`$ is key: it allows for condition (46) to be satisfied with finite $`t`$. A pulse $`f(t)`$ has a continuous band of frequencies of width $`\mathrm{\Gamma }1/\tau _g`$ (where $`\tau _g`$ is the gate duration), with amplitudes (Fourier transform) $`\widehat{f}(\mathrm{\Omega })`$, which add to and smear the effective Bohr spectrum $`\{\omega \}`$. If the pulse is long (a slow gate) then only a narrow band appears, and the smearing effect is unimportant. More precisely, if $`1/\tau _g`$ is much smaller than the typical difference of the Bohr frequencies, the โenergy quantaโ $`m\mathrm{\Omega }`$ \[with $`|m|`$ restricted by the (typically small) dimension of the system Hilbert space\] cannot fill the gap between $`\omega `$ and $`\omega ^{}`$ and the condition (46) can be satisfied. This is our adiabatic approximation. For fast pulses, when $`1/\tau _g`$ is comparable to $`|\omega \omega ^{}|`$, the condition (46) cannot be fulfilled: the effective Bohr spectrum becomes quasi-continuous and the denominator in condition (46) becomes abitrarily small. The result is that the WCL analysis breaks down and non-Markovian effects dominate.
Thus, the condition for the adiabatic limit (Markov approximation valid) is: โthe width of the band is much smaller than the minimal difference of the effective Bohr frequenciesโ. This is in contradiction with the fast gate assumption, A1.
### III.3 Section Summary
The main advantage of the MME (23) is its consistency with thermodynamics. Namely, as a consequence of the KMS condition (15) and the condition (18), for a generic initial state the system tends to its thermal equilibrium (Gibbs) state at the temperature of the heat bath Alicki:87 . (An important exception to this rule are states within a decoherence-free subspace Zanardi:97c ; LidarWhaley:03 , but these states are not generic due to required symmetry properties of the system-bath interaction.) Therefore the dissipative part of the generator must depend strongly on the Hamiltonian dynamics. This is consistent with the notion of a coarse-grained description familiar from the study of MMEs: the bath needs a time much longer than $`\mathrm{max}_{\omega _{kl}}1/\omega _{kl}`$ to โlearnโ the systemโs Hamiltonian in order to drive it to a proper Gibbs state. In other words, *the Markov approximation is, equivalently, a long-time limit* (compared to $`\mathrm{max}_{\omega _{kl}}1/\omega _{kl}`$ โ the systemโs Bohr frequencies), and one cannot expect this approximation to be valid at short times. However, FT-QEC assumes operations on a time-scale that is short on the scale set by $`\mathrm{max}_{\omega _{kl}}1/\omega _{kl}`$.
Strictly speaking the MME (23) is valid only when $`H_S`$ is not time dependent. As we have shown, one can relax this by assuming slowly varying $`H_S`$, giving rise to an โadiabatic MMEโ, Eqs. (35), (36) and (25). However, to accept Eqs. (35), (36) and (25) as a genuine Markovian description is somewhat of a stretch, since the real question is not whether one obtains the Lindblad form, but rather *how* $`(t)`$ *depends on the Hamiltonians* $`H_S(u)`$*, locally (i.e.* $`ut`$*) or nonlocally. For fast gates and generic environments the dependence is non-local, involving memory effects.* In any case, the crucial condition that must be satisfied for a (generalized) MME is Eq. (46), which implies that the average Bohr spectrum must be discrete. In essence, as long as the applied control does not spoil this discreteness a (generalized) MME can be derived. On the other hand, *this means that fast gates are incompatible with the MME*, in violation of A1 of FT-QEC theory. The corollary: *finite speed of gates implies non-Markovian effects*.
## IV Are the Standard FT-QEC Assumptions Internally Consistent?
We now briefly summarize our examination of the assumptions of FT-QEC, in light of the considerations above, and highlight where there may be internal inconsistencies in FT-QEC. As discussed above, there are essentially two rigorous approaches to the derivation of the MME: (i) the SCL, which is compatible with arbitrarily fast Hamiltonian manipulations, but requires the high-$`T`$ limit; (ii) the WCL, which is compatible with thermodynamics at arbitrary $`T`$, but requires adiabatic Hamiltonian manipulations.
The standard theory of FT-QEC (excluding Refs. Terhal:04 ; Aliferis:05 ; Aharonov:05 ) requires a quantum computer (QC) undergoing Markovian dynamics, supplemented with a constant supply of cold and fresh ancillas. These assumptions are contradictory under the SCL, since the QC would have to be at high-$`T`$, while the ancillas require low-$`T`$ on the same energy scale $`E`$ (set by the Bohr energies of the system = computer \+ ancillas). Specifically, if we were to assume that for the ancillas too $`kTE`$, they would quickly become highly mixed. If we insist that $`kTE`$ for the ancillas, then by coupling them to the QC we can no longer assume, in the SCL, that the total system = QC + ancillas is described by Markovian dynamics.
If, on the other hand, we approach the problem from the (physically more consistent) WCL, then A3 is incompatible with A1 (the assumption of fast gates). Namely, in the WCL only adiabatic Hamiltonian manipulations are allowed. Specifically, the Markov approximation in the WCL requires a *discrete* system (effective) Bohr frequency spectrum, such that the condition$`\tau _g\mathrm{max}_{\omega _{kl}}1/\omega _{kl}`$ can be satisfied, hence violating the $`\tau _g\omega _B=O(\pi )`$ condition of A1. These conclusions are unavoidable if one accepts thermodynamics, since they follow from seeking a Markovian master equation that satisfies the KMS condition โ a necessary condition for return to thermodynamic equilibrium in the absence of external driving. We take here the reasonable position that a fault tolerant QC cannot be in violation of thermodynamics.
## V Possible objections to the inconsistency
In this section we analyze a list of possible objections to the inconsistency we have pointed out.
### V.1 Is thermodynamics relevant?
With respect to the SCL: โ*Thermodynamics is irrelevant (since a QC need not ever be in thermal equilibrium)*.โ
Note that we never claim that the QC is in thermal equilibrium; only the bath is. This assumption is a simplification which allows us to use a single parameter $`T`$ and therefore a single โthermal memory timeโ $`\mathrm{}/kT`$. There is no reason to use a nonthermal bath or many heat baths with different temperatures: this does not make the spectral density flat and can only introduce more parameters.
### V.2 Doesnโt the interaction picture save the day?
With respect to the WCL: โ*Suppose we have the following Hamiltonian in the Schrodinger picture:* $`H=H_S+H_C(t)+H_{SR}+H_R`$ *where* $`H_SH_C=`$*control Hamiltonian* $`H_{SR}`$*. Then in the interaction picture with respect to* $`H_S`$ *the term* $`H_C`$ *is dominant and hence can implement fast gates. However, in the original Schrรถdinger picture* $`H_C`$ *is small and hence the adiabatic limit for the derivation of the MME is satisfied. Thus we have an example where we can have fast gates (in the interaction picture) and still the WCL can be satisfied so that the Markovian limit can be reached. Moreover, this is the relevant limit relevant for quantum optics, e.g., trapped ions.*
There are a number of problems with this argument. First, one should be more careful about the formulation of the condition for adiabaticity. It can be stated as $`|d\omega (t)/dt|\omega (t)^2`$, where $`\omega (t)`$ is a โrelevantโ Bohr frequency. Merely comparing norms as above does not guarantee adiabaticity. Second, in the quantum optics context we note the following. For three-level trapped ions we have two Bohr frequencies: a large, time-independent $`\omega _1`$, and a small, time-dependent $`\omega _2(t)`$ (degenerate levels splitting). Only $`\omega _2`$ is โrelevantโ because it is related to gates, and then the adiabatic condition implies that $`|d\omega _2(t)/dt|`$ is correspondingly small, which contradicts the fast gate condition A1. Third, the inequality $`H_CH_{SR}`$ is in fact not satisfied in the Markovian WCL, where $`H_{SR}`$ diverges (one should not confuse the small system-reservoir coupling parameter involved in the van-Hove limit with the operator norm, which can be infinite).
### V.3 Doesnโt quantum optics provide a counterexample?
With respect to the WCL: โ*Trapped ions and other quantum optics systems provide a counter-example: a system experimentally satisfying Markovian dynamics and allowing fast Rabi operations*.โ
We have already addressed quantum optical systems in Section III.2.2. Let us add a few comments. We do not know of any quantum optics experiment testing the Markov approximation with the accuracy relevant for FT-QEC (for quantum dots, on the other hand, non-Markovian effects are very visible). We know that for constant, and also for strictly periodic Hamiltonians (which corresponds in quantum optics to a constant external laser field), the Davies derivation can be applied (or extended, as in Section III.2) and the Markov approximation is applicable. The problem appears for fast gates. It would be difficult to test the Markov approximation in this case with the required accuracy, because, e.g., the results depend on the shape of the pulse. A relevant example is resonance fluorescence, as described in Carmichael:book \[pp.43-61\], and as discussed in Section III.2.2. The damping effects are only present in the widths of spectral lines โ see Carmichael:book \[p.61, Fig. 2.5\]. The Markov approximation gives Lorentzians while non-Markovian dynamics may give rise to more complicated lineshapes. Consider a 2-level atom like in Carmichael:book Section 2.3.2., and in particular the final formula Eq. (2.96), which describes resonance fluorescence via a MME. The author claims that for typical parameters in quantum optics the dissipative part does not depend on the Rabi frequency $`\mathrm{\Omega }_\mathrm{R}`$ \[recall our discussion in Section III.2.2\]. Hence, as the gates are entirely related to $`\mathrm{\Omega }_\mathrm{R}`$, it appears that either fast or slow gates are possible. The argument is based on the small ratio $`\mathrm{\Omega }_\mathrm{R}/\omega _A<10^{10}/10^{15}`$ (where $`\omega _A`$ is the Bohr frequency). This is fine for replacing the spectral density at $`\omega _A\pm \mathrm{\Omega }_\mathrm{R}`$ by the density at $`\omega _A`$, but the subsequent argument that we can replace \[in Eq. (2.94)\] $`\mathrm{\Omega }_\mathrm{R}`$ by $`0`$ is inaccurate. This would be correct only if the decay time $`\tau =1/\gamma `$ is short enough such that $`\mathrm{\Omega }_\mathrm{R}\tau 1`$. However, as explained in Section III.2.2, in this case the Davies type averaging makes no sense physically. In fact, typically for radiation damping $`\tau =10^8`$s, and then $`\mathrm{\Omega }_\mathrm{R}\tau <100`$ only. Hence for a fixed $`\mathrm{\Omega }_\mathrm{R}`$ we do in fact not have a simple Lindblad generator (of the type (2.96) in Carmichael:book ), but rather a more complicated generator with Lindblad operators depending on the Rabi frequency, as in Eq. (47). Again, in the derivation of a proper generator an averaging over terms of the form $`\mathrm{exp}(i\mathrm{\Omega }_\mathrm{R}t)`$ must be performed. Therefore the condition for the adiabatic approximation involves the Rabi frequency $`\mathrm{\Omega }_\mathrm{R}`$ and cannot be satisfied for fast gates. For experiments based on *spectral measurements* the difference between the two types of generators we have just discussed is probably irrelevant for many reasons; however, the quantum state of the atom at a given moment is sensitive to a small change in the Lindblad operators, and this is important in a fault tolerant implementation of quantum logic gates.
### V.4 Is A1 truly an assumption of FT-QEC?
With respect to the WCL: โ*Doesnโt A1 impose an unnecessary constraint on FT-QEC, in that gates are not required to satisfy the condition* $`\tau _g\omega =O(\pi )`$*?*
In other words, one might argue in favor of slow gates, where instead the condition is $`\tau _g\omega O(\pi )`$. Such gates are certainly relevant in the context of the adiabatic quantum computing (AQC) paradigm Farhi:01 , holonomic QC ZanardiRasetti:99 ; ZanardiRasetti:2000 , or topological quantum computing (TQC) Kitaev:97 ; Freedman:01 ; DasSarma:05 . We comment in more detail on AQC, HQC, and TQC in Section VI. The question of interest to us is whether an adiabatic gate satisfying $`\tau _g\omega O(\pi )`$ is applicable to the standard FT-QEC paradigm we are considering here, and which is very different from AQC, HQC, and TQC.
First, let us clarify that by gates we mean one and two-qubit unitaries picked from well-known discrete and small sets of universal gates Nielsen:book . An algorithm is constructed via a sequence of such gates, and computational complexity is measured in terms of the minimal number of required gates. Of course one can instead join all gates used in a given algorithm into a single unitary and call this a gate, but then one runs into the problem of finding a relevant (physical) Hamiltonian and quantifying computational complexity. For a given gate there are infinitely many Hamiltonian realizations. Among these are fast ones (optimal) which satisfy $`\tau _g\omega =O(\pi )`$ and slow ones (adiabatic) satisfying $`\tau _g\omega O(\pi )`$ (all inequalities here are in the sense of orders of magnitude). For example, consider a $`\pi `$-rotation. The fast (optimal) realization satisfies $`\tau _g\omega =\pi `$ (compatible with A1), while the slow (adiabatic) one satisfies $`\tau _g\omega =\pi +2\pi n`$ with $`n1`$ (contradicts A1).
Now, one may ask whether a slow realization of gates can prevent the inconsistency with the WCL. We argue, based on computational complexity considerations, that the answer to this question is negative. To see this, note first that non-Markovian errors are uncorrectable in standard FT-QEC. Therefore such non-Markovian, uncorrectable errors accumulate during the computation (by definition, they are not corrected by โMarkovian FT-QECโ), and in order to keep them under control, the probability of such errors per gate, $`p_{\mathrm{non}\mathrm{M}}`$, should scale as
$`p_{\mathrm{non}\mathrm{M}}`$ $``$ $`O[1/(\text{volume of algorithm})]`$ (51)
$`=`$ $`O[1/(\text{input size})^\alpha ],`$
where $`\alpha `$ is some fixed power. Now, it follows from our discussion in Section III.2 that the more adiabatic the evolution, the smaller is the probability of the non-Markovian errors per gate. Therefore, if one writes the adiabaticity condition as $`\tau _g\omega >M`$, where $`M1`$ is the โadiabatic slowness parameterโ, then the probability of non-Markovian errors should satisfy
$$p_{\mathrm{non}\mathrm{M}}O(1/M^\beta ),$$
(52)
where $`\beta `$ is another fixed power \[$`\omega `$ (the Bohr or Rabi frequency) is limited essentially by the choice of physical system\]. Comparing the two expressions for $`p_{\mathrm{non}\mathrm{M}}`$, we see that $`M`$ must grow with input size. This means that if one works with adiabatic gates in order to keep the dynamics (approximately) Markovian, the result is that one must slow the gates in proportion to the input size (to some power). This, however, violates the threshold condition of FT-QEC, in which the input size and gate times are independent parameters (see, e.g., Theorem 12 in Aharonov:99 ).
### V.5 Measurements
With respect to both the WCL and the SCL: โ*Recent results on fault-tolerant QC using measurements only (e.g., Nielsen:04 ; Raussendorf:05 ) render all the claimed problems irrelevant*.โ
Indeed, we have so far discussed only the problems with quantum logic gates. Moreover, measurements are an integral part of FT-QEC theory as well, in particular to reset and disentangle ancillas before they are introduced into an error-correction circuit. Therefore some remarks on the use of measurements are in order.
In the most advanced FT-QEC scheme of Aharonov:99 , measurements are performed at the end of the computation. However, this approach demands a high resource overhead, which may make it impractical. Therefore, more recent proposals (e.g., Knill:05 ; Steane:04 ) rely on feedback mechanisms employing the results of quantum measurements. Those โmeasurements in the middle of computationโ are treated for simplicity as certain von-Neumann projective measurements (but with efficiency $`1`$) satisfying a *repeatability condition*. The latter implies that the subsequent measurements reduce the measurement error exponentially as their number increases. This assumption should be carefully scrutinized, within realistic Hamiltonian models of quantum measurement treated as a dynamical process. Here, again one can expect that the tacit assumption of statistical independence of repeated measurements is in conflict with the non-Markovian character of the dynamics of open quantum systems.
As all proposed measurement schemes are based on electromagnetic interactions, it should be possible to construct a rather general Hamiltonian framework and apply it to various particular implementations. Indeed, this has been done, e.g., for a single-electron tunneling (SET) transistor coupled capacitively to a Josephson junction qubit Shnirman:98 . Rather than assuming that the measurement apparatus is coupled to the system whenever measurements must be performed โ an option which is hard to achieve in mesoscopic systems โ Ref. Shnirman:98 makes the reasonable assumption that the measurement apparatus is always coupled to the system, but is in a state of equilibrium when it is not needed. A measurement is then performed by driving the measuring device out of equilibrium, in a manner that dephases the qubit to be measured. Generic features emerging from this analysis are the existence of three different time-scales characterizing the measurement: the dephasing time, the measurement time (which may be longer than the dephasing time), and the mixing time (the time after which all the information about the initial quantum state is lost due to the transitions induced by the measurement). Ref. Shnirman:98 thus arrives at a criterion for a โgoodโ quantum measurement: the mixing time should be longer than the measurement time. A time-scale analysis of measurements in optical systems, accounting for spontaneous emission, can be found, e.g., in Ref. Teich:89 . A fully consistent analysis of FT-QEC should account for the existence of such time-scales in a dynamic description of the measurement process. In particular, it is important to set appropriate bounds on these time-scales, so that they may be taken into account in a threshold calculation (an analysis based on a stochastic error model was reported in Ref. Steane:03 ).
### V.6 Degenerate Qubits
With respect to the SCL: โ*Degenerate qubits automatically satisfy the high* $`T`$ *limit since their intrinsic energy scale vanishes*.โ
Examples of degenerate qubits are common, e.g., in trapped ion quantum computing implementations where a pair of degenerate hyperfine states can serve as a qubit, with an auxiliary third level used to implement quantum logic gates via Raman transitions Wineland:98 . The case of degenerate qubits is somewhat more subtle to analyze within the context we have explained above. Naively, in such a case the high-$`T`$ limit is indeed automatically satisfied, since the system energy scale is zero. Therefore it appears that one could claim that the SCL version of the Markov approximation is attainable. However, upon closer examination this still seems problematic. Indeed, the vanishing of an energy scale for degenerate qubits holds, strictly speaking, only for fully adiabatic techniques, e.g., HQC ZanardiRasetti:99 ; ZanardiRasetti:2000 . Otherwise transformations between logical states are achieved by resorting to effective Hamiltonians which involve *virtual* transitions. For instance, if $`|0`$ and $`|1`$ denote degenerate qubit levels (e.g., hyperfine levels of an ion), one can introduce far-detuned (e.g., laser) couplings of $`|0`$ and $`|1`$ with a third auxiliary level. Second order perturbation theory then yields the effective Hamiltonian $`H_{\mathrm{eff}}=(\mathrm{\Omega }_\mathrm{R}^2/\mathrm{\Delta })|10|+\mathrm{h}.\mathrm{c}.`$, where $`\mathrm{\Omega }_\mathrm{R}`$ and $`\mathrm{\Delta }`$ are the laser Rabi coupling and detuning, respectively. Therefore we see that an effective, small but non-vanishing, energy scale $`E_1:=\mathrm{\Omega }_\mathrm{R}^2/\mathrm{\Delta }`$ is introduced. (Note that in order for perturbation theory to be valid one must have $`\mathrm{\Omega }_\mathrm{R}\mathrm{\Delta }`$, which in turn implies $`E_1\mathrm{\Delta }`$.) Yet another energy scale is provided by the spectral width $`E_2`$ of the laser pulse shape $`\mathrm{\Omega }_\mathrm{R}(t)`$; in order to suppress unwanted *real* transitions, one must impose in addition that $`E_2\mathrm{\Delta }`$. At any rate, the appearance of these new system-energy scales implies that once again the SCL-type contradiction applies. On the other hand, we can make both $`E_1`$ and $`E_2`$ small at the price of lengthening the gating time ($`\tau _g\mathrm{max}\{1/E_1,\mathrm{\hspace{0.17em}1}/E_2\}`$). This implies, once again, an adiabatic limit and the applicability of the WCL. Therefore it appears that as long as one restricts manipulations to adiabatic ones (thus contradicting A1), quantum computing with degenerate qubits is possible even in the Markovian limit. We expand on this viewpoint below.
### V.7 Impure Ancillas
With respect to the SCL: โ*Do ancillas really need to be pure?*
What precisely is the role of the ancillas in QEC? A popular answer is that they serve as an โentropy sinkโ for the errors accumulated during the quantum computation. This entropy in the system arises from the entanglement between system and bath, and the role of the ancillas is to remove this entanglement. I.e., in a perfect quantum error correction step the entanglement between system and bath is transferred to the ancillas and bath. A natural objection to our SCL-based inconsistency is to claim that, in fact, ancillas need not be pure, or could perhaps even be highly mixed. However, this is not supported by the (current) standard theory of FT-QEC. Consider, e.g., an error correction circuit based on the Steane 7-qubit code. It takes as input ancillas prepared in the $`|\psi _a=(|0_L+|1_L)/\sqrt{2}`$ state, where $`|0_L`$ and $`|1_L`$ are codewords. The physical qubits which comprise such ancillas, are coupled bitwise via CNOT gates to the physical qubits making up the encoded data qubits in the circuit. If instead we input an ancilla in a mixed state, this is equivalent to inputting a classical mixture with erred codewords, e.g., $`(1p)|\psi _a\psi |+p|\varphi _a\varphi |`$, where $`|\varphi _a`$ is an erred codeword. If one of these errors is a phase-flip, it feeds back (via the CNOT gates) into the data qubits, producing an error Gottesman:97a . Without fault-tolerance this means that there are now two errors (in the ancillas block and the data block), which may be more than the code can handle. In FT-QEC theory such errors are accounted for, but their magnitude is bounded from above (e.g., $`p`$ in the above example must be small). We note that an ancilla which is initially entangled with the data qubits (violating the assumption of being introduced into the circuit in a tensor-product state) is essentially equivalent to the case of an impure ancilla just described (tracing over the data qubits yields an impure ancilla state).
A more general approach showing the importance of the assumption of pure ancillas is the following (fairly standard account of QEC).
i) Preparation.โ
Let the initial state of system + reservoir + ancillas, with respective Hilbert spaces $`_S,_R,_A`$, be: $`\rho _{SRA}^0=|\psi _S\psi _S||0_R0_R|\rho _A`$, where we have allowed for ancillas in a mixed state $`\rho _A`$.
ii) System-reservoir interaction (decoherence).โ
$$\rho _{SRA}^0\stackrel{U_{SR}}{}\rho _{SRA}^1=\underset{e,e^{}}{}U_e|\psi _S\psi _S|U_e^{}|e_Re_R^{}|\rho _A,$$
(53)
where $`e`$โs denote the *errors* belonging to the set $``$ that the code $`๐`$ can correct, and where $`|e_R`$ are the corresponding states of the reservoir. The error operators $`U_e`$ are assumed to be unitary and with linear span of dimension $`||`$.
iii) System-ancilla interaction (syndrome extraction).โ
This interaction takes the form $`U_{SA}=_e\mathrm{\Pi }_eT_e`$ where the $`T_e`$โs are unitaries over $`_A`$ such that $`T_e|0_A=|e_A`$ and $`\mathrm{\Pi }_eI_๐|ee|`$ <sup>13</sup><sup>13</sup>13We know that $`_S๐๐ฎ๐`$ \[$`๐ฎ`$=syndrome subsystem, dim$`๐=||`$; $`๐`$=remainder (=$`0`$ for subspace-based codes)\] Knill:97b ; Knill:99a ; Zanardi:99d .:
$`\rho _{SRA}^1`$ $`\stackrel{U_{SA}}{}`$ $`\rho _{SRA}^2`$
$`=`$ $`{\displaystyle \underset{e,e^{}}{}}U_e|\psi _S\psi _S|U_e^{}^{}|e_Re_R^{}|T_e\rho _AT_e^{}.`$
iv) Error recovery.โ
Unitary recovery is implemented via $`\stackrel{~}{U}_{SA}=||^{1/2}_eU_e^{}I_R|e_Ae_A|`$, where for unitarity we need $`e_A|e_A^{}=\delta _{e,e^{}}`$. By applying $`\stackrel{~}{U}_{SA}`$ and tracing over both $`R`$ and $`A`$ (assuming the $`|e_R`$โs too are orthonormal) one obtains
$$\rho _S^{\mathrm{out}}=\frac{1}{||}\underset{e,fE}{}U_f^{}U_e|\psi _S\psi _S|U_e^{}U_ff_A|T_e\rho _AT_e^{}|f_A.$$
(55)
In the case of a pure ancillas $`\rho _A=|0_A0_A|`$ one has $`f_A|T_e\rho _AT_e^{}|f_A=|f_A|e_A|^2=\delta _{f,e}`$ and therefore the ideal case $`\rho _A^{\mathrm{out}}=|\psi _S\psi _S|`$ is recovered. One can also consider the fidelity
$`F`$ $`:=`$ $`\psi _S|\rho _S^{\mathrm{out}}|\psi _S`$
$`=`$ $`||^1{\displaystyle \underset{e,fE}{}}|\psi _S|U_f^{}U_e|\psi _S|^2f_A|T_e\rho _AT_e^{}|f_A.`$
Provided the error operators $`U_f`$ satisfy the condition for a non-degenerate code $`\psi _S|U_f^{}U_e|\psi _S=\delta _{f,e}`$ Knill:97b , one obtains $`F=||^1_ee_A|T_e\rho _AT_e^{}|e_A=0_A|\rho _A|0_A.`$ Clearly, $`F=1`$ iff $`\rho _A=|0_A0_A|,`$ i.e., *the ancillas are pure*. One can also consider non-unitary recovery via ancilla measurements and conditional unitaries, with Kraus operators given by $`A_e=||^{1/2}U_e^{}I_R|e_Ae_A|`$. The conclusion that the ancillasโ state must be pure is unchanged.
We note that FT is obtained by adding concatenation and, in steps iii) and iv), preparing and coupling encoded ancillas with the system in a suitable way, e.g., as in the Steane-code example above. In this case it is permissible to allow slightly impure ancillas, and relax the assumptions that, in step ii) the environment couples only to the system, and in steps iii,iv), the environment does not act. This formulation, however, does not allow arbitrarily mixed-state ancillas, as argued in the Steane-code example. While such a formulation of FT-QEC theory might still emerge (for example, by using algorithmic cooling techniques Schulman:98 ; Schulman:05 , which, however, at present assume perfect gates), it does not appear possible at present to relax the assumption of cold ancillas.
### V.8 Hot QC, cold ancillas, and fast QC-ancilla interactions in the SCL
With respect to the SCL: *โOne can keep the ancillas coupled to a separate cold bath and then couple them for only a short time to the QC: what matters then is the* $`\text{T}_1`$ *timescale and that one can be very long compared to the required ancilla-QC coupling timeโ*.
Let us paraphrase this objection. If one can make $`H_{SA}`$ (system-ancillas) very large then one could beat the rate of ancilla heating by strongly coupling the QC and ancillas. I.e., suppose one would like to bring the ancillas in from their cold reservoir to couple to the system, which is coupled to a hot reservoir as required for the SCL. The ancillas then heat up fast, but there is a timescale associated with this heating (โ$`T_1`$โ), which one wishes to beat. Now if one could make the system-ancilla coupling very strong then one could, presumably, use the ancillas (e.g. for syndrome extraction) faster than their heating rate, while they are still sufficiently pure for fault tolerance purposes.
The simplest argument against this objection is the following. In the setting of the objection, the QC is described by the SCL (high $`T`$) while the ancillas are described by the WCL (low $`T`$). Strong and fast coupling between the QC and the ancillas is unacceptable according to the WCL because it is fast (only adiabatic manipulations are allowed), and according to the SCL because it is strong (โstrongโ refers to the systemโs Hamiltonian part, while in the SCL this Hamiltonian is weak in comparison with the system-bath coupling).
However, one could go on to argue that the ancillas are a different species than the QC qubits, and in particular have a different intrinsic (less dense) energy scale, so that they are at low $`T`$ on the scale set by the QC qubits. In this case both ancillas and QC are described by the SCL. Then the problem with the objection is the following: recall that in the SCL (see Section III.1) one must rescale $`H_{SR}`$ and $`H_{AR}`$ as $`H_{SR}/ฯต`$ and $`H_{AR}/ฯต`$ respectively, where here $`R`$ denotes the common reservoir the system and the ancillas are coupled to. The heating rate is proportional to the square of the coupling strength to the reservoir, i.e., to $`1/ฯต^2`$, and hence diverges in the SCL. Therefore to beat the ancilla heating process via fast manipulation of the system-ancilla coupling one would have to rescale $`H_{SA}`$ at least by $`1/ฯต^2`$, but this contradicts the SCL derivation, where in fact one must keep $`H_{SA}`$ fixed while rescaling $`H_{SR}`$. The reason for this is that, in the SCL derivation, it is the system (now including the ancillas) that sets the timescale against which reservoir correlations must be accelerated.<sup>14</sup><sup>14</sup>14Let us also consider the issue from the perspective of thermodynamics. This is not really necessary, since the arguments above about the SCL are rigorous, but is interesting in its own right. First, we remark that error correction should really be made to work at the common lower (initial ancillasโ) temperature. Heating a part of a QC only to be closer to the Markovian limit is a suboptimal strategy, because it increases the strength of the noise and stimulates entropy production. Second, in standard FT-QEC heat (entropy) flows from the QC to the ancillas only, while in reality one should expect a flow in both directions and additionally an entropy production. To see this let us ignore for the moment the coupling of the QC to the bath, and consider ancillas coupled to a heat bath at temperature $`T`$. The ancillas can be kept pure by maintaining an energy gap $``$ $`kT`$. Assume that the initial state of QC ($`C`$) and ancillas ($`A`$) is a product state $`|\psi _C|\psi _A`$. Switching on the interaction $`H_{CA}`$ we induce an equilibration process (because the dynamics is Markovian) of $`C+A`$ towards the Gibbs state $`\rho _{CA}=\mathrm{exp}(H_{CA}/kT)/Z`$, which is *entangled* (here for simplicity $`H_{CA}`$ contains not only the interaction but is the total Hamiltonian of $`C+A`$). After a single step of error correction the total state of $`C+A`$ can be modeled by $`(1p)U|\psi _C|\psi _AU^{}+p\rho _{CA}`$, where $`U`$ is unitary and $`0<p1`$. Then we switch off the interaction with the ancillas. Whatever we do next separately with $`C`$ and $`A`$, we cannot eliminate the error due to the entanglement present in the term $`p\rho _{CA}`$. This type of incorrectable error accumulates and destroys FT-QEC. This is the back flow of entropy from the ancillas bath to the QC, mentioned above.
## VI Alternatives to Markovian FT-QEC
### VI.1 Nature of the non-Markovian errors in the WCL
While we have pointed out that, in the WCL, the application of fast gates is likely to violate the conditions required for Markovian dynamics to persist, we have not been specific about the type of non-Markovian effects that will emerge. It is well known that FT-QEC is capable of dealing with errors that change due to the application of gates. Namely, assume (slightly) faulty gates correcting a specific error model described by a CP map $`\mathrm{\Lambda }`$ \[recall Eq. (1)\], are applied in sequence, $`\mathrm{\Lambda }U_N^{}\mathrm{\Lambda }U_{N1}^{}\mathrm{}\mathrm{\Lambda }U_1^{}`$, and these gates are (in some appropriate norm) close to the ideal gates $`\{U_i\}_{i=1}^N`$, i.e., for all $`i`$, $`U_i^{}U_i^{}I1`$. Then by inserting $`U_i^{}U_i`$โs everywhere one obtains the new sequence $`\mathrm{\Lambda }_NU_N\mathrm{\Lambda }_{N1}U_{N1}\mathrm{}\mathrm{\Lambda }_1U_1`$, where $`\mathrm{\Lambda }_i:=\mathrm{\Lambda }U_i^{}U_i^{}`$, and FT-QEC is capable of dealing with such a (gate-modified) error model. However, the non-Markovian effects that arise due to the application of fast gates in the WCL, will in general *not* be describable by a simple time-local modification such as $`\mathrm{\Lambda }\mathrm{\Lambda }U_i^{}U_i^{}`$. This can be worked out, e.g., using the methods of Ref. Breuer:02 .
In order to formulate consistent alternatives to standard, Markovian FT-QEC theory, it seems useful to start with a Hamiltonian formulation. As the discussion below will illustrate, it appears that a hybrid approach will be necessary, which combines alternatives to standard QC with a new version of FT-QEC.
### VI.2 Adiabatic Quantum Computing (AQC)
We keep A2 and A3, discard A1 (fast gates), and work in a purely adiabatic mode, thus permitting a consistent WCL. This may indeed be possible using the adiabatic quantum computing (AQC) approach of Farhi et al. Farhi:01 . At present there is little understanding of the fault-tolerance of AQC. Some recent works explore AQC in the presence of decoherence and/or control errors Childs:02 ; Shenvi:03 ; Roland:04 ; Aberg:04 ; SarandyLidar:05 ; aberg-2005-72 . Indeed, the subject of the adiabatic approximation in open quantum systems has only very recently been addressed SarandyLidar:04 , and used to study AQC in open systems SarandyLidar:05 . Error correcting codes for AQC were introduced very recently in Jordan:05 , but the corresponding universal Hamiltonians involve many-body interactions (four and six-body for 1-local and 2-local errors, respectively).
### VI.3 Holonomic Quantum Computing (HQC)
Another possibility for keeping A2 and A3, and discarding A1, is provided by HQC ZanardiRasetti:99 ; ZanardiRasetti:2000 . HQC is an adiabatic scheme which relies on Abelian or non-Abelian geometric phases to implement quantum logic gates. Quantum information is encoded in a *degenerate* set of eigenstates of a Hamiltonian depending on a set of controllable parameters, e.g., external laser fields (recall our discussion of degenerate states above). When these are adiabatically driven along a suitable closed path, the initial quantum state is transformed by a non-trivial unitary transformation (holonomy) that is geometrical in nature. The key point is that the geometrical nature of the quantum holonomies is believed to render HQC inherently robust against certain kinds of errors. This alleged fault-tolerance has only recently been seriously begun to be examined Solinas:04 ; the emerging picture is that, while stability against decoherence must still be assessed, HQC seems to exhibit a strong robustness against stochastic errors in the control process generating the required adiabatic loops Zhu:04 . Moreover, in the adiabatic limit of Markovian dynamics it has been show that the geometric phase of a single qubit coupled to a magnetic field is robust against both dephasing and spontaneous emission (but not against bit flips) SarandyLidar:05a . Nevertheless, since deviations from strict adiabaticity are inevitable, and adiabaticity is particularly challenging to satisfy in open quantum systems SarandyLidar:04 , it is tempting to combine HQC with FT-QEC in order to address the performance of HQC in the presence of decoherence errors. Alternatively, we note that a hybrid approach that seems to be rather promising is the embedding of HQC within a DFS WuZanardiLidar:05 . This amounts to realizing a set of universal quantum gates, acting on a DFS, by means of non-abelian quantum holonomies. This strategy brings together the โbest of two worldsโ: the quantum decoherence avoidance virtues of DFSs and the fault-tolerance of the all-geometric holonomic control. It is possible that such an approach can be implemented for quantum information processing in, e.g., trapped ions and quantum dots.
### VI.4 Topological Quantum Computing (TQC)
A robust way of performing quantum computations is based on excitations with fractional statistics, since they have several fault-tolerant properties built in. This idea is known as topological quantum computing (TQC) Kitaev:97 ; Freedman:01 ; DasSarma:05 . Physical realizations of the simplest versions of TQC have been considered in the literature, using, e.g., rotating Bose-Einstein condensates Paredes:01 and superconducting circuits Ioffe:02 . Let $`๐`$ denote the manifold of quantum codewords. Strikingly, in TQC, one can have a trivial Hamiltonian, e.g., $`H|_๐=0`$, but nevertheless obtain non-trivial quantum evolution due to the existence of an underlying topological *global* structure (boundary conditions). Quantum encoding is typically performed in a properly designed degenerate ground state $`๐`$. This fact implies, for low enough temperature, an exponential suppression of errors on encoded quantum information due to thermal fluctuations. More importantly, topological features can render such a ground state stable against errors represented by *local* operators, namely error operators that do not involve a number of qubits of the order of the size of the system. For example, in the so-called toric codes Kitaev:97 ; Kitaev:book , qubits are encoded in the ground-state manifold of a lattice of interacting spins in such a way that degenerate ground states are mutually connected only via high powers (scaling linearly with lattice size) of local operators. Thus, here the fault-tolerance properties are already built-in at the *physical level*. However, while one can argue that topological encoding provides a stable and passive quantum memory, it is not self-correcting as in todayโs โeffectively naturally fault-tolerantโ classical architectures (see Ref. Bacon:05 for an eloquent exposition of this point). Moreover, it is important to realize that as far as we know, in its present state TQC still requires active intervention, in the form of FT-QEC, when one tries to compute fault-tolerantly. Indeed, Preskill writes in Ref. Preskill:TQC-notes \[p.62\], โIt is therefore implicit that the temperature is small enough compared to the energy gap of the model that thermally excited anyons are too rare to cause trouble, that the anyons are kept far enough apart from one another that uncontrolled exchange of charge can be neglected, and in general that errors in the topological quantum computation are unimportant. If the error rate is small but not completely negligible, then the standard theory of quantum fault tolerance can be invoked to boost the accuracy of the simulation as neededโ. Ref. Dennis:02 takes this approach and explicitly lists A2 and A3 as necessary requirements for fault-tolerant TQC. In contrast, A1 is definitely *not* required in TQC: one performs computations by adiabatically dragging quasiparticles around one another, and these operations must be slow relative to the gap between the ground state and the first excited state. The larger the gap the easier it is to satisfy this adiabaticity condition, so this requirements is compatible with the thermal suppression of errors mentioned above. In addition, TQC requires time-dependent controls to read out the encoded data (Ref. Freedman:00a shows that all measurements can be postponed until the readout of the final result of the computation). However, a fully Hamiltonian analysis of the fault-tolerance of such measurements is still lacking. Nevertheless, one could argue that the error rate in a topological quantum computer could be made arbitrarily small by increasing the system size and careful engineering, so that (similarly to todayโs self-correcting, fault-tolerant classical computers), one could ultimately perform TQC without any active intervention other than read-out of the encoded data. An interesting, recent development in this direction was reported in Ref. Bacon:05 , which suggests that certain three-dimensional quantum spin-lattices might be self-correcting.
### VI.5 Non-Markovian Quantum Computing
We keep A1 and A2 but discard A3 (the Markov approximation) at least in part. This appears to be a reasonable approach in many cases, since the Markov approximation is clearly a highly idealized limit (though it does hold remarkably well in some optical systems and in liquid state NMR). Indeed, the degree of accuracy to which the Markov approximation must be satisfied has been quantified, e.g., by Steane in Steane:04 : the probability of an uncorrectable (i.e., non-Markovian) error per gate must be $`<`$ $`10^{10}`$ for a computation involving $`10^9`$ gates (this probability must scale with the input size, as explained in Section V.4). Alternative approaches to dealing with non-Markovian baths are therefore of interest. For example, the papers Terhal:04 ; Aliferis:05 ; Aharonov:05 present an extension of FT-QEC theory to a non-Markovian setting. We offer in this context the following observations:
1. An important ingredient carried over directly and without change from Markovian FT-QEC theory, is the crucial role of the fresh and nearly pure ancillas Terhal:04 ; Aliferis:05 ; Aharonov:05 . We believe that the detailed mechanism for introducing and discarding ancillas at specific times should be reconsidered within a fully Hamiltonian framework.
2. As recognized and discussed in Terhal:04 , the important assumption of a small norm of the system-bath interaction Hamiltonian (e.g., Eq. (58) in Ref. Aliferis:05 ) is not satisfied for some standard models of open systems. For example, a linear coupling to a bosonic heat bath involves unbounded interaction Hamiltonians and a high-frequency cutoff. In general, the assumption of a small norm of the system-bath interaction Hamiltonian is much stricter than the WCL and is not satisfied for most standard models of reservoirs.
Another approach to fault-tolerance in a non-Markovian setting is the recently developed time-concatenated dynamical decoupling method KhodjastehLidar:04 (see also Viola:02 for a version of dynamical decoupling with bounded-strength controls). However, comment 2. above about the small norm of the system-bath interaction Hamiltonian applies here as well. Therefore more general methods are required to deal with the full scope of baths one can expect in quantum computing implementations. A promising possibility in this direction is to incorporate fault-tolerant dynamical decoupling in a feedback loop.
## VII Conclusions
We have listed a set of minimal assumptions made in the theory of fault-tolerant quantum error correction (FT-QEC): 1) fast gates (on the timescale set by the inverse of the relevant Bohr or Rabi frequency), 2) a supply of fresh and nearly pure ancillas, 3) a Markovian bath.
We have also reviewed the only two known rigorous general limits leading to Markovian dynamics: the singular coupling limit (SCL), which involves taking a high temperature limit, and the weak coupling limit (WCL), which requires either a constant or an adiabatic system Hamiltonian, and averaging over long times in comparison with the inverse of the relevant Bohr frequency. These two limits allow one to replace the reservoir autocorrelation function by a Dirac delta, which leads to the Markovian limit.
A close examination of the assumptions of FT-QEC has led us to conclude that assumption 3 can be sustained together with assumption 1 in the SCL, and together with assumption 2 in the WCL. However, it is not possible to maintain all three assumptions in either the SCL or the WCL. We therefore conclude that, at present, there exists an inconsistency in the formulation of the theory of FT-QEC for Markovian baths. We have also listed a number of alternatives to Markovian FT-QEC which, from the point of view adopted here, are free of inconsistencies. However, none of these alternatives is so comprehensive as to include the full range of errors one might expect in a full-scale implementation of quantum computing. In particular, recent results on fault tolerance in non-Markovian settings Terhal:04 ; Aliferis:05 ; Aharonov:05 ; KhodjastehLidar:04 , while representing a significant step forward, make a crucial assumption about the smallness of the norm of the system-bath interaction Hamiltonian, which severely restricts the class of physical reservoirs.
###### Acknowledgements.
We thank Dave Bacon, Andrew Doherty, Daniel Gottesman, Hideo Mabuchi, John Preskill, Alireza Shabani, and Barbara Terhal for very useful discussions (though this does not imply their agreement with our conclusions). Their insightful comments helped us sharpen our critique and formulate the questions in Section V. R.A. thanks for the support from the Polish Ministry of Science and Information Technology- grant PBZ-MIN-008/P03/2003 and the EC grant RESQ IST-2001-37559, D.A.L. thanks the Sloan Foundation for a Research Fellowship and the DARPA-QuIST program for support. P.Z. acknowledges support by the European Union FET project TOPQIP (Contract No. IST-2001-39215).
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# hep-th/0506083 Exploring Vacuum Manifold of Open String Field Theory
## 1 Introduction
String field theory (SFT) is expected to provide a framework where distinct string backgrounds can be studied in terms of some universal set of underlying degrees of freedom. For example, it was conjectured that SFT admits the closed string vacuum solution describing the state after all unstable D-branes are allowed to decay via condensation of the tachyon . If it happens, no open string excitations appear around the tachyon vacuum, and the energy difference between the originally defined D-brane vacuum and the tachyon vacuum is precisely equal to the brane tension .
It is highly desirable to find analytic candidates for the classical solution describing the tachyon vacuum. They will make it possible not only to justify the above conjecture but to give valuable insights into the structure of the vacuum state in SFT. Such a candidate was discussed in Refs. . There, given was a series of one parameter families of classical solutions associated with the functions $`h_a^l(w)`$. The solutions were constructed with these functions which specify combinations of the BRS current and the ghost field operated on the identity state. They are labeled by an integer $`l`$, and classified in such a way that non-trivial solutions emerge for the function $`h_{a_b}^l(w)`$ with the boundary value $`a=a_b`$, while all other solutions associated with $`h_a^l(w)(aa_b)`$ become gauge transformations of the D-brane vacuum. In addition to these, we have constructed another class of solutions with higher order zeros .
Based on the analysis on the cohomology , scattering amplitudes and the potential height around each solution, we understand that all of the above non-trivial solutions are equally valid candidates for the tachyon vacuum. This makes us to believe that all of them are actually equivalent and related by appropriate symmetry transformations present in the SFT action. To clarify this point is the main motivation of this paper.
The SFT action has the gauge symmetry and it is also invariant under transformations generated by particular combinations of the Virasoro operators, $`K_n=L_n()^nL_n`$ . Later we will see that the symmetry generated by $`K_n`$ belongs to the โglobalโ part of the gauge symmetry. The solutions, if they are really equivalent, are to be related by some gauge transformations. In this paper, we find operators that transform the higher $`l`$ solutions down to the $`l=1`$ solution. This is the main result of this paper. The operators are written in terms of generators $`K_n`$ with even $`n`$. One may wonder what would happen to the inverse of the relations. When we go back to a higher $`l`$ solution from a lower $`l`$ solution, we encounter a subtlety, that will be explained in a later section and further discussion will be found in the last section.
This paper is organized as follows. In the next section, we briefly review how analytic classical solutions to the SFT action are constructed in Refs. . The action expanded around a solution carries a new BRS charge in its kinetic term. The charge seems to have a peculiar ghost structure. However, it will be explained that this is naturally understood from the first quantization point of view. In section 3, we discuss the symmetries of the SFT action. In particular, the deformation of a classical solution under $`K_n`$-transformations are described. Since the interaction vertex is invariant under the transformation, the change of classical solution has an effect only on the kinetic term, or the BRS charge, in the expanded action. Therefore, when we have actions expanded around two classical solutions that are related by this type of symmetry, we ought to be able to find an appropriate transformation between two BRS charges. This is the subject of the section 4. We construct the operator $`U`$ relating the BRS charges for the $`l>1`$ and $`l=1`$ classical solutions and some properties of the operator are reported. The operator has a well-defined normal ordered expression and generate a sound string field transformation on the component fields. The last section is devoted to discussions. We have added four appendices. The appendix A is to evaluate the action for the string field configuration in a pure gauge form. Some technical points in relation to section 3 are explained in appendix B. In appendix C, matrix elements of $`U`$ are calculated. The universal solutions due to are not in the Siegel gauge as shown in appendix D.
## 2 Classical Solutions in String Field Theory
This section is a short summary of how analytic classical solutions are constructed and what makes us believe that they really correspond to the tachyon vacuum. Though the construction of solutions for SFT looks quite non-trivial, the first quantization point of view provides us with an intuitive picture of the solutions: it also helps us to understand the constraint embedded in the new BRS charge, defined on a non-trivial solution. The first quantization point of view will be explained in the second subsection.
### 2.1 Classical solutions and gauge transformations
We summarize our construction of the classical solutions, paying special attention to gauge transformations in cubic string field theory (CSFT).
The action in the original CSFT ,<sup>1</sup><sup>1</sup>1For later convenience, we write the BRS charge dependence of the action as $`S[\mathrm{\Psi },Q_B]`$.
$`S[\mathrm{\Psi },Q_B]={\displaystyle \frac{1}{g^2}}{\displaystyle \left(\frac{1}{2}\mathrm{\Psi }Q_\mathrm{B}\mathrm{\Psi }+\frac{1}{3}\mathrm{\Psi }\mathrm{\Psi }\mathrm{\Psi }\right)},`$ (2.1)
is characterized by the BRS charge $`Q_\mathrm{B}`$ and the $``$-product, which enjoy the following properties:
1. $`Q_BA=0`$;
2. $`Q_B(AB)=(Q_BA)B+()^{|A|}AQ_BB`$;
3. $`(AB)C=A(BC)`$;
4. $`AB=()^{|A||B|}BA`$.
By $`|A|`$, we denote the Grassmannian even-oddness of the string field $`A`$: $`|A|=+1(1)`$, when A is Grassmann even (odd).
The action (2.1) is an analogue to the integration of the Chern-Simons three form : the string field $`\mathrm{\Psi }`$, the integration $``$ and the BRS charge $`Q_\mathrm{B}`$ correspond to the connection $`A`$, the integration of differential forms $``$ and the exterior derivative $`d`$ of ordinary differential geometry. Similarly to the Chern-Simons action, it is easy to show that the action (2.1) transforms as
$`S[\mathrm{\Psi }^{},Q_B]=S[\mathrm{\Psi },Q_B]+S[g^1Q_\mathrm{B}g,Q_B]`$ (2.2)
under the gauge transformation
$`\mathrm{\Psi }^{}=g^1Q_\mathrm{B}g+g^1\mathrm{\Psi }g.`$ (2.3)
Here, the string functional $`g`$ is an element of the stringy gauge group in which the multiplication law is given by the star product. This stringy gauge group is expected to have much richer structure than that of the ordinary Yang-Mills theory. Eq. (2.2) reminds us of a similar expression for the Chern-Simon theory (based on a compact group). Though it is tempting to think of a concept of homotopy class for this stringy gauge group, we do not have much more to discuss along this direction. For the purpose of the present paper, it is suffice to know that the second term in eq. (2.2) vanishes for the functional $`g`$ connected to the identity $`I`$ via a continuous deformation (see appendix A for the proof). In other words, the action is invariant under such a gauge transformation.
The equation of motion is given by the variational principle,
$`Q_\mathrm{B}\mathrm{\Psi }+\mathrm{\Psi }\mathrm{\Psi }=0.`$ (2.4)
The lhs of (2.4) is in the form of the โfield strengthโ. Therefore, at least formally, โflat connectionsโ are classical solutions to string field theory. Let us write such a classical solution as
$`\mathrm{\Psi }_0=g_0^1Q_\mathrm{B}g_0,`$ (2.5)
where $`g_0`$ is a (group-valued) string functional.
If we expand the string field around the classical solution (2.5) as
$`\mathrm{\Psi }=g_0^1Q_\mathrm{B}g_0+\mathrm{\Phi },`$ (2.6)
the action (2.1) becomes
$`S[\mathrm{\Psi },Q_B]=S[g_0^1Q_\mathrm{B}g_0,Q_B]{\displaystyle \frac{1}{g^2}}{\displaystyle \left(\frac{1}{2}\mathrm{\Phi }Q_\mathrm{B}^{}\mathrm{\Phi }+\frac{1}{3}\mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi }\right)},`$ (2.7)
where the new BRS charge $`Q_\mathrm{B}^{}`$ is defined as
$`Q_\mathrm{B}^{}AQ_\mathrm{B}A+g_0^1Q_\mathrm{B}g_0A(1)^{\left|A\right|}Ag_0^1Q_\mathrm{B}g_0,`$ (2.8)
for an arbitrary string field $`A`$.
After having considered the general structure of the string field theory action, let us explicitly describe the classical solutions given in Refs..
Consider the specific gauge functional
$`g_0(h)`$ $`=`$ $`\mathrm{exp}(q_\mathrm{L}(h)I)`$ (2.9)
$`=`$ $`Iq_\mathrm{L}(h)I+{\displaystyle \frac{1}{2!}}q_\mathrm{L}(h)Iq_\mathrm{L}(h)I+\mathrm{},`$
where the operator $`q_\mathrm{L}`$ is defined in terms of the ghost number current $`J_{\mathrm{gh}}(w)`$ and a function $`h(w)`$ satisfying $`h(\pm i)=0`$ and $`h(1/w)=h(w)`$:
$`q_\mathrm{L}(h)={\displaystyle _{C_{\mathrm{left}}}}{\displaystyle \frac{dw}{2\pi i}}h(w)J_{\mathrm{gh}}(w).`$ (2.10)
The integration path indicated by the subscript $`C_{\mathrm{left}}`$ is over the left half of the string, ie, $`\pi /2<\sigma <\pi /2`$ on the unit circle for the variable, $`w=e^{i\sigma }`$. The gauge functional (2.9) gives rise to a classical solution $`\mathrm{\Psi }_0(h)=g_0(h)^1Q_\mathrm{B}g_0(h)`$, which may be rewritten as
$`\mathrm{\Psi }_0(h)=Q_\mathrm{L}(e^h1)IC_\mathrm{L}((h)^2e^h)I,`$ (2.11)
where the operators $`Q_\mathrm{L}`$ and $`C_\mathrm{L}`$ are defined with the BRS current and the ghost field:
$`Q_\mathrm{L}(f)={\displaystyle _{C_{\mathrm{left}}}}{\displaystyle \frac{dw}{2\pi i}}f(w)J_\mathrm{B}(w),C_\mathrm{L}(f)={\displaystyle _{C_{\mathrm{left}}}}{\displaystyle \frac{dw}{2\pi i}}f(w)c(w).`$ (2.12)
The solution (2.11) has a well-defined Fock space expression in the universal subspace, spanned by the matter Virasoro generators and ghost oscillators acting on the $`SL(2,R)`$ invariant vacuum. Therefore, it would be appropriate to call the solution (2.11) as the universal solution.
For the solution (2.11), the new BRS charge can be expressed as
$`Q_\mathrm{B}^{}=Q(e^h)C((h)^2e^h).`$ (2.13)
The operators $`Q`$ and $`C`$ are defined as
$`Q(f)={\displaystyle \frac{dw}{2\pi i}f(w)J_\mathrm{B}(w)},C(f)={\displaystyle \frac{dw}{2\pi i}f(w)c(w)}.`$ (2.14)
Here the integrations are over the unit circle.
Formally, the solution (2.11) is in a pure gauge form, and, therefore, could be gauged away. As shown in , however, this is not always the case: some non-trivial solutions emerge at the boundary of one parameter deformation of a certain class of functions chosen for $`h`$. The functions are given by
$`h_a^l(w)=\mathrm{log}\left\{1{\displaystyle \frac{a}{2}}(1)^l\left(w^l\left({\displaystyle \frac{1}{w}}\right)^l\right)^2\right\}(l=1,2,3,\mathrm{}).`$ (2.15)
We have a series of functions labeled by the integer $`l`$. Accordingly, we have a series of non-trivial solutions associated with the functions. The solutions will be addressed as TTK solutions in this paper.
Now we describe evidences suggesting that the TTK solutions are really non-trivial. The reality of the function, for $`w`$ on the unit circle, restricts the parameter $`a`$ to be $`a1/2`$. This condition guarantees the hermiticity of the corresponding classical solution as well as the new BRS charge, as expected from eq. (2.13). Since $`h_{a=0}^l(w)=0`$ and the corresponding classical solution vanishes, it would be reasonable to expect that the solution written with $`h_a^l`$ is a trivial pure gauge for sufficiently small $`a`$. Indeed, we know the following facts for $`a>1/2`$:
1. The expanded action can be transformed back to the action (2.1) ;
2. The new BRS charge provides us with the same cohomology as the original BRS charge ;
3. The expanded theory reproduces ordinary open string scattering amplitudes ;
4. A numerical study with the level truncation technique shows that the expanded theory has a non-perturbative vacuum and its vacuum energy tends to the D-brane tension as the truncation level increases .
These facts are consistent with our expectation that the solution is trivial pure gauge. However, at the boundary $`a=1/2`$, the expanded theory shows completely different properties<sup>2</sup><sup>2</sup>2It is not known whether the non-trivial solutions with $`a=1/2`$ can still be written in the form $`g^1Q_\mathrm{B}g`$. If it is the case, we may consider that the solutions are obtained by โlarge gauge transformations.โ That would be a strong evidence for some topological structure of the stringy gauge group.:
1. The new BRS charge has the vanishing cohomology in the Hilbert space with the ghost number one ;
2. The vanishing of open string scattering amplitudes, the result is consistent with the absence of open string excitations (no open string theorem) ;
3. We can show numerically that the non-perturbative vacuum found for $`a>1/2`$ disappears as the parameter $`a`$ approaches to $`1/2`$ .
The above results implies that the solution with $`a=1/2`$ indeed corresponds to the tachyon vacuum.
### 2.2 Interpretation of new BRS charges in the first quantized theory
The emergence of the non-trivial theory for $`a=1/2`$ can also be seen from the first quantization point of view.
In the original action (2.1), we have the Kato-Ogawaโs BRS charge
$`Q_\mathrm{B}={\displaystyle \frac{dw}{2\pi i}\left[c(w)T_X(w)+(bcc)(w)\right]},`$ (2.16)
where $`T_X(w)`$ is the stress tensor for the string coordinates $`X`$ and $`b`$ is the antighost. The BRS charge $`Q_\mathrm{B}`$ is known to be constructed directly from the first-class constraint $`T_X(w)0`$ .
We now consider a modification of the constraint surface by multiplying a function, $`e^{h(w)}`$. Then, the modified BRS charge constructed from the constraint $`e^{h(w)}T_X(w)0`$ takes of the form
$`Q_\mathrm{B}^{}={\displaystyle \frac{dw}{2\pi i}e^{h(w)}\left[c(w)T_X(w)+(bcc)(w)+\frac{1}{2}c(w)\{(h)^2+3(^2h)\}\right]}.`$ (2.17)
Here the term linear in the ghost is needed to ensure the nilpotency condition $`(Q_\mathrm{B}^{})^2=0`$. The expression (2.17) coincides with (2.13) where the BRS current is given by $`j_B(w)=c(w)[T_X(w)+(bc)(w)+3/2^2c(w)]`$. It means that the replacement of the BRS charge $`Q_\mathrm{B}`$ by $`Q_\mathrm{B}^{}`$ corresponds to the replacement of the constraint $`T_X(w)0`$ by $`e^{h(w)}T_X(w)0`$. This change of the constraint can be absorbed by a redefinition of ghost and antighost:
$`c(w)c(w)e^{h(w)}`$ $`=`$ $`e^{q(h)}c(w)e^{q(h)},`$
$`b(w)b(w)e^{h(w)}`$ $`=`$ $`e^{q(h)}b(w)e^{q(h)},`$ (2.18)
so that
$`e^{q(h)}Q_\mathrm{B}e^{q(h)}=Q_\mathrm{B}^{}.`$ (2.19)
This relation holds, of course, only if the operator $`e^{q(h)}`$ with
$`q(h)={\displaystyle \frac{dw}{2\pi i}h(w)J_{\mathrm{gh}}(w)}`$ (2.20)
is well-defined. The operator $`e^{q(h_a^l)}`$, with the function given in (2.15), is well-defined for $`a>1/2`$, but not for $`a=1/2`$ . Whether the similarity transformation (2.19) makes sense or not depends on the distribution of zeros of the function $`\mathrm{exp}(h_a^l(w))`$: all zeros are distributed off the the unit circle $`|w|=1`$ for $`a>1/2`$, while they merge on the unit circle for $`a=1/2`$. This change in the distribution of zeros may be related to the non-trivial modification of the constraint in first quantized theory: the constraint surface given by $`\mathrm{exp}(h_{1/2}^l(w))T_X(w)0`$ becomes physically distinct from the original surface $`T_X(w)\mathrm{exp}(h_{a>1/2}^l(w))T_X(w)0`$.
In this section, the properties of our classical solutions are summarized. As we have seen, there present an infinite number of non-trivial solutions and various results suggest that they all describe the tachyon vacuum. If it is really the case, they are equivalent with each other and related via the symmetries of CSFT. The symmetries of CSFT are the subject of the next section.
## 3 Symmetries of CSFT and Classical Solutions
The CSFT action has a subalgebra of the Virasoro algebra as its symmetry . Here we will see that this symmetry may be considered as a subgroup of the โglobalโ part of the stringy gauge symmetry. Therefore it provides a way to relate classical solutions and, thereby the SFT actions defined on them.
The subalgebra is generated by $`K_n=L_n()^nL_n`$, where $`L_n`$ is the Virasoro operator. Using the properties ,
$`{\displaystyle K_nA}`$ $`=`$ $`0\mathrm{for}\mathrm{all}A,`$
$`K_n(AB)`$ $`=`$ $`(K_nA)B+AK_nB,`$ (3.1)
it is easy to show the invariance of the action under an infinitesimal transformation generated by $`K_n`$ :
$$\delta SK_n\left(\frac{1}{2}\mathrm{\Psi }Q_B\mathrm{\Psi }+\frac{1}{3}\mathrm{\Psi }\mathrm{\Psi }\mathrm{\Psi }\right)=0.$$
(3.2)
We now consider a particular type of gauge transformation and find it to be a finite form of $`K_n`$-transformation written as
$$\mathrm{\Psi }^{}e^{K(v)}\mathrm{\Psi },$$
(3.3)
with $`K(v)_{n>0}v_nK_n`$. The parameters $`v_n`$ will be specified below.
Let us take $`u_L(f)\mathrm{exp}\left(๐ฏ_L(f)I\right)`$ as a gauge functional. The operator $`๐ฏ_L(f)`$ is defined as
$$๐ฏ_\mathrm{L}(f)=_{C_{\mathrm{left}}}\frac{dw}{2\pi i}f(w)T(w),$$
(3.4)
similarly to eq. (2.12), with the total energy-momentum tensor $`T(w)`$. It is easy to see that the gauge transformation (2.3) with $`u_L(f)`$ can be written as
$`\mathrm{\Psi }^{}`$ $`=`$ $`u_L^1\mathrm{\Psi }u_L=U(f)\mathrm{\Psi },`$ (3.6)
$`U(f)\mathrm{exp}\left({\displaystyle \frac{dw}{2\pi i}f(w)T(w)}\right).`$
Note that the first term in (2.3) with the BRS charge is absent in the above expression since $`[Q_B,๐ฏ_L(f)]=0`$ and $`Q_BI=0`$. In deriving the last expression of eq. (3.6), we used the properties of the half splitting operators in eqs. (B.6) and (B.8) that require the function $`f(w)`$ to satisfy the condition, $`f(w)=(dw/d\stackrel{~}{w})f(\stackrel{~}{w})`$ for $`\stackrel{~}{w}=1/w`$.<sup>3</sup><sup>3</sup>3Some more detailed derivation of (3.6) is described in appendix B. Expanding the function as $`f(w)=_nv_nw^{n+1}`$, we find the relation, $`v_n=v_n()^{n+1}`$, from the condition. Using this relation, it is easy to see that the integral in eq. (3.6) becomes the operator $`K(v)_{n>0}v_nK_n`$.
The absence of the term $`u_L(f)Q_Bu_L(f)^1`$ in the expression (3.6) allows us an interesting interpretation of the transformation (3.3). The BRS charge $`Q_B`$ corresponds to the external derivative $`d`$ in the Chern-Simons theory. A gauge transformation in the CS-theory with the absence of the derivative term is simply a global transformation. Similarly, a stringy gauge transformation without the first term in (2.3) may be considered to be a โstringy global transformationโ in SFT. Note that, strictly speaking, a global transformation is not necessarily a finite $`K_n`$-transformation. The type of global transformations written as (3.3) forms a specific subgroup in the stringy global transformations. In the rest of the paper, we discuss only this subgroup of the global symmetry.
Now, we apply the global transformations discussed above on a given classical solution. Let us find the action for the fluctuation $`\mathrm{\Phi }`$ around a classical solution $`\mathrm{\Psi }_0`$. Substituting $`\mathrm{\Psi }\mathrm{\Psi }_0+\mathrm{\Phi }`$ into the action $`S[\mathrm{\Psi },Q_B]`$, we obtain
$`S[\mathrm{\Psi }_0+\mathrm{\Phi },Q_B]`$ $`=`$ $`S[\mathrm{\Psi }_0,Q_B]+S[\mathrm{\Phi },Q_B(\mathrm{\Psi }_0)].`$ (3.7)
Here $`Q_B(\mathrm{\Psi }_0)`$ is defined as
$`Q_B(\mathrm{\Psi }_0)A`$ $``$ $`Q_BA+\mathrm{\Psi }_0A()^{|A|}A\mathrm{\Psi }_0`$ (3.8)
on a string field $`A`$. The nilpotency follows from the equation of motion for $`\mathrm{\Psi }_0`$.
Once we find a classical solution, we can, at least formally, obtain other solutions related by the gauge symmetry. As for solutions related as eq. (3.3), it is easy to see the following statements to hold:
1. If $`\mathrm{\Psi }_0`$ solves the string equation of motion, ie, $`Q_B\mathrm{\Psi }_0+\mathrm{\Psi }_0\mathrm{\Psi }_0=0`$, then $`\mathrm{\Psi }_0^{}e^{K(v)}\mathrm{\Psi }_0`$ is also a solution;
2. Furthermore the BRS charges defined around two solutions are related as
$$e^{K(v)}Q_B(e^{K(v)}\mathrm{\Psi }_0)e^{K(v)}=Q_B(\mathrm{\Psi }_0).$$
(3.9)
Since the transformation (3.3) leaves the action invariant, the first statement is trivial. Technically speaking, it can be shown by using the property,
$$(e^{K(v)}A)(e^{K(v)}B)=e^{K(v)}(AB),$$
(3.10)
which is obtained from eq. (3.1). The second statement follows from the relation,
$`Q_B(\mathrm{\Psi }_0^{})e^{K(v)}A=e^{K(v)}Q_B(\mathrm{\Psi }_0)A,`$ (3.11)
on a generic string field $`A`$. Eq. (3.11) may be derived from the definition of the BRS charge given in eq. (3.8).
Now we would like to see how a universal solution may be deformed by the action of the $`K_n`$ operators. Leaving discussions of the finite transformation in the next section, here we consider the change of a solution under an infinitesimal transformation and discuss its implication.
The generic form of the classical solutions obtained in Refs. is given as
$`|\mathrm{\Psi }_0`$ $`=`$ $`Q_L(F)|I+C_L(G)|I,`$ (3.13)
$`G(w)={\displaystyle \frac{(F(w))^2}{1+F(w)}},`$
where $`F(z)`$ is an analytic function satisfying the relations $`F(1/w)=F(w)`$ and $`F(\pm i)=0`$.
It is easy to see that the action of $`e^{K(\epsilon )}`$ on the universal solution $`|\mathrm{\Psi }_0`$ produces yet another universal solution, at least in the first order in deformation parameters $`\epsilon _n`$. The effect of the operator appears as a change in the function $`F(z)`$,
$`|\stackrel{~}{\mathrm{\Psi }}_0=e^{K(\epsilon )}|\mathrm{\Psi }_0`$ $``$ $`Q_L(\stackrel{~}{F})|I+C_L(\stackrel{~}{G})|I,`$ (3.14)
$`\stackrel{~}{F}(w)`$ $``$ $`F(w){\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\epsilon _nu_n(w)F(w),`$ (3.16)
$`u_n(w)w(w^n()^nw^n).`$
It is easily confirmed that the deformed function also satisfy two conditions for a classical solution.<sup>4</sup><sup>4</sup>4In eq. (3.16), we observe that no choice for a set of parameters leaves the function invariant. This implies that the symmetry generated by $`K_n`$ does not survive, at least, in its original form.
The above calculation shows that we may explore the submanifold of classical solutions by the action of $`e^{K(v)}`$. This opens a possibility to relate different solutions written in the same form as described in (3.13). In particular, the series of solution constructed in may be related by the operator $`e^{K(v)}`$ with its parameters appropriately chosen. Of course, this cannot be realized by infinitesimal transformations discussed here and we have to consider finite transformations.
We make another observation that supports this idea. Consider the SFT defined around the classical solution $`\mathrm{\Psi }_0(h_a^l)`$ written as eq. (2.11) with the function in (2.15). After taking the limit of $`a\frac{1}{2}`$, we have the action $`S[\mathrm{\Phi },Q_B^{(l)}]`$ with the BRS charge given as
$`Q_B^{(l)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}Q_B+{\displaystyle \frac{()^l}{4}}(Q_{2l}+Q_{2l})+2l^2\left(c_0{\displaystyle \frac{()^l}{2}}(c_{2l}+c_{2l})\right)`$ (3.17)
$`=`$ $`Q(F^{(l)})+C(G^{(l)}),`$
where the moments of the BRS current are defined in the expansion, $`J_B(w)_nQ_nw^{n1}`$, and the functions $`F^{(l)}(w)`$ and $`G^{(l)}(w)`$ are given as
$`F^{(l)}(w)`$ $``$ $`{\displaystyle \frac{()^l}{4}}\left(w^l+(w)^l\right)^2,G^{(l)}(w)l^2w^2()^l\left(w^l(w)^l\right)^2.`$ (3.18)
If the classical solutions labeled by $`l`$ and $`m(lm)`$ are really related by the finite form of the $`K_n`$ symmetry, the BRS charges in the actions $`S[\mathrm{\Phi },Q_B^l]`$ and $`S[\mathrm{\Phi },Q_B^m]`$ are to be related as
$`Q_B^{(m)}=e^{K(v)}Q_\mathrm{B}^{(l)}e^{K(v)},`$ (3.19)
with some operator $`e^{K(v)}`$, as we stated in (3.9). From eq. (3.17) and the commutation relation, $`[K_n,Q_{2l}]=2l(Q_{2l+n}()^nQ_{2ln})`$, we realize that $`K_n`$ with $`n=`$ even are to be used for the purpose. So it is possible for eq. (3.19) to hold. We are to find out whether we may choose proper parameters so that eq. (3.19) holds. Further discussion of the transformation (3.19) will be given in the next section.
Before closing this section, we introduce operators which play an important role in relating BRS charges:
$$U_{2l}(t)=\mathrm{exp}\left[\frac{()^l}{4l}\mathrm{ln}\left(\frac{1t}{1+t}\right)K_{2l}\right](l=1,2,\mathrm{}).$$
(3.20)
When ordered with respect to the Virasoro operators, they are expressed as
$$U_{2l}(t)=\mathrm{exp}\left(\frac{()^lt}{2l}L_{2l}\right)\mathrm{exp}\left(\frac{1}{2l}\mathrm{ln}(1t^2)L_0\right)\mathrm{exp}\left(\frac{()^lt}{2l}L_{2l}\right).$$
(3.21)
The operators in eq. (3.20) generate the conformal transformations ,
$$f_{2l}(z,t)=\left(\frac{z^{2l}()^lt}{1()^ltz^{2l}}\right)^{\frac{1}{2l}}.$$
(3.22)
A comment is in order. These conformal transformations are used earlier in connection with the TTK solutions. The parameter $`t`$ introduced here corresponds to $`Z(a)(1+a\sqrt{1+2a})/a`$ in the earlier expression.
## 4 Relating classical solutions with different values of $`l`$
When two classical solutions are related by the global transformation, expansions around them produce BRS charges satisfying eq. (3.9). Here we present a way to construct operators relating the BRS charges for the TTK classical solutions.
In the last section, we have seen that the BRS charges may be related as eq. (3.19) with the operator $`K(v)`$ written in terms of $`K_n`$ ($`n=`$even). Here we construct the operator $`e^{K(v)}`$ that relates charges associated with $`l=1`$ and $`l1`$. In the next subsection, it will be shown that the operator is realized in the limit of a one-parameter family of operators, $`U(t)(1t0)`$:
$`U(t=0)=1,Q_B^{(l=1)}=\underset{t1}{lim}U(t)Q_B^{(l)}U^1(t).`$ (4.1)
Unfortunately, there seems to be a subtlety in this operator: when trying to obtain the higher $`l`$ BRS charge from the $`l=1`$ charge, we encounter a problem. This will be explained briefly in the second subsection. It is not clear to us at this moment whether this is just a technical problem or much deeper one. In the third subsection, the properties of the operator $`U(t)`$ are investigated. In particular, it will be shown that it has the well-defined normal ordered expression in terms of the Virasoro generators even in the limit of $`t1`$.
### 4.1 Higher $`l`$ solutions down to $`l=1`$ solution
Our construction of operators that relate BRS charges is based on the observation to be explained below. On the $`l`$-th BRS charge (3.17), we act the operator introduced in eq. (3.20) and find the charge transformed as,
$`U_{2l}(t)Q_B^{(l)}U_{2l}^1(t)=Q(F_1^{(l)})+C(G_1^{(l)}),`$ (4.2)
$`F_1^{(l)}(w,t)`$ $``$ $`F^{(l)}(z)|_{z=f_{2l}(w,t)}={\displaystyle \frac{()^l}{4}}\left(w^l+(w)^l\right)^2{\displaystyle \frac{(1+t)^2}{\left(1+()^ltw^{2l}\right)\left(1+()^ltw^{2l}\right)}}`$ (4.3)
$`=`$ $`{\displaystyle \frac{1+t}{2}}+{\displaystyle \frac{1t^2}{4}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}()^{n1}()^{ln}t^{n1}(w^{2ln}+w^{2ln}),`$
$`G_1^{(l)}(w,t)`$ $``$ $`G^{(l)}(z)|_{z=f_{2l}(w,t)}\times \left({\displaystyle \frac{df_{2l}(w,t)}{dw}}\right)^2`$ (4.4)
$`=`$ $`()^ll^2w^2\left(w^l(w)^l\right)^2{\displaystyle \frac{(1+t)^2(1t)^4}{\left(1+()^ltw^{2l}\right)^3\left(1+()^ltw^{2l}\right)^3}}`$
$`=`$ $`{\displaystyle \frac{2l^2(1+t+t^2)}{w^2(1t^2)}}{\displaystyle \frac{l^2}{2w^2(1t^2)}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}()^{ln}g_n(t)(w^{2ln}+w^{2ln}).`$
Here $`f_{2l}(w,t)`$ is given in eq. (3.22) and the coefficients $`g_n(t)`$ are given by
$`g_n(t)=(t)^{n1}\left[n^2+n+4(n+1)t2(n^22)t^24(n1)t^3+(n^2n)t^4\right].`$ (4.5)
Substituting eqs. (4.3) and (4.4) into the terms on the rhs of eq. (4.2), we obtain the expression of $`U_{2l}(t)Q_B^{(l)}U_{2l}^1(t)`$ as
$`U_{2l}(t)Q_B^{(l)}U_{2l}^1(t)={\displaystyle \frac{1+t}{2}}Q_B`$ $`+`$ $`{\displaystyle \frac{1t^2}{4}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(1)^{n(l1)1}t^{n1}(Q_{2ln}+Q_{2ln})`$
$`+{\displaystyle \frac{2(1+t+t^2)l^2}{1t^2}}c_0`$ $``$ $`{\displaystyle \frac{l^2}{2(1t^2)}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}()^{nl}g_n(t)(c_{2nl}+c_{2ln}).`$ (4.6)
On the rhs, we have higher moments of the BRS currents and pure ghost terms.
Note that, in the limit of $`t1`$, all the terms containing the moments of BRS currents vanish. The remaining pure ghost terms become divergent. So some care should be taken when we take this limit at this stage. It would be worth pointing out that the pure ghost terms are indeed those appeared in the vacuum string field theory. This fact was first realized in Refs. and utilized recently to regularize the VSFT .
Another important observation on eq. (4.6) is the fact that the rhs contains the BRS charge itself. Since all the TTK solutions have this property, we may relate various BRS charges via the expression on the rhs of (4.6). In concrete, we act $`U_2(t)`$ for $`l=1`$ on the rhs of (4.6) and see what would come out. Our expectation is that the result is somewhat close to the BRS charge for the $`l=1`$ solution.
The action of $`U_2(t)`$, on the rhs of (4.6), replaces the argument of $`F_1^{(l)}(w,t)`$ by $`f_2(w,t)`$. After some calculations, we obtain
$`F_2^{(l)}(w,t)`$ $``$ $`F_1^{(l)}(z,t)|_{z=f_2(w,t)}`$ (4.7)
$`=`$ $`{\displaystyle \frac{1}{4}}\left(w1/w\right)^2\left({\displaystyle \frac{2}{l+1}}\right)^2{\displaystyle \frac{1}{\left(1+\frac{l1}{l+1}w^2\right)\left(1+\frac{l1}{l+1}w^2\right)}}+O\left((1+t)\right).`$
Similarly we obtain the expression for $`G_2^{(l)}(w,t)`$. Note that the singular behavior in the limit of $`t1`$, observed in the rhs expression of eq. (4.6), is absent in (4.7). The action of $`U_2(t)`$ has canceled the singular behavior. In the limit both $`F_2^{(l)}(w,t)`$ and $`G_2^{(l)}(w,t)`$ are finite. So we may take the limit of $`t1`$ in the expression for $`U_2(t)^1U_{2l}(t)Q_B^{(l)}U_{2l}(t)^1U_2(t)`$ where $`U(t)U_2(t)^1U_{2l}(t)`$:
$`\underset{t1}{lim}U_2(t)^1U_{2l}(t)Q_B^{(l)}U_{2l}(t)^1U_2(t)Q(\stackrel{~}{F}_2^{(l)})+C(\stackrel{~}{G}_2^{(l)}),`$ (4.8)
$`\stackrel{~}{F}_2^{(l)}(w)\underset{t1}{lim}F_2^{(l)}(w,t)={\displaystyle \frac{1}{4}}\left(w1/w\right)^2\left({\displaystyle \frac{2}{l+1}}\right)^2{\displaystyle \frac{1}{\left(1+\frac{l1}{l+1}w^2\right)\left(1+\frac{l1}{l+1}w^2\right)}},`$ (4.9)
$`\stackrel{~}{G}_2^{(l)}(w){\displaystyle \frac{(\stackrel{~}{F}_2^{(l)}(w))^2}{\stackrel{~}{F}_2^{(l)}(w)}}.`$
Here we find the function $`(w1/w)^2/4`$ in eq. (4.9), the function for the $`l=1`$ BRS charge. The remaining factor appeared in eq. (4.9) may be removed by an appropriate transformation. Indeed, we can deform the $`l=1`$ BRS charge to the form of eq. (4.9):
$`U_2\left({\displaystyle \frac{l1}{l+1}}\right)Q_B^{(1)}U_2\left({\displaystyle \frac{l1}{l+1}}\right)^1=Q(\stackrel{~}{F}_2^{(l)})+C(\stackrel{~}{G}_2^{(l)}).`$ (4.10)
Combining eqs. (4.8) and (4.10), we finally reach the $`l=1`$ BRS charge starting from the higher $`l`$ charge:
$`U_2^1\left({\displaystyle \frac{l1}{l+1}}\right)\underset{t1}{lim}\left(U_2(t)^1U_{2l}(t)Q_B^{(l)}U_{2l}(t)^1U_2(t)\right)U_2\left({\displaystyle \frac{l1}{l+1}}\right)=Q_B^{(1)}.`$
The expression is also rewritten as
$`\underset{t1}{lim}U(t)Q_B^{(l)}U^1(t)=Q_B^{(1)},`$ (4.11)
where $`U(t)U_2(\stackrel{~}{t})^1U_{2l}(t)`$ with $`\stackrel{~}{t}`$ defined as
$`\stackrel{~}{t}{\displaystyle \frac{(l1)+(l+1)t}{(l+1)(l1)t}}.`$ (4.12)
The parameter in $`U_2`$ is now $`\stackrel{~}{t}`$ due to the extra action of $`U_2(\frac{l1}{l+1})`$ in eq. (4.10).
### 4.2 $`l=1`$ solution up to $`l=2`$ solution
A natural question is whether we can obtain the SFT action around the $`l=2`$ solution starting from that for $`l=1`$ solution. As a relation of BRS charges, our question may be formulated as follows: is it possible to find an operator $`๐ฐ`$ such that $`๐ฐ^1Q_B^{(1)}๐ฐ=Q_B^{(2)}`$? In this subsection, we explain what we have understood in relation to this question.
Let us reconsider eq. (4.11) for $`l=2`$. Before taking the limit, we write
$`Q_B^{(1)}(t)U_2^1(\stackrel{~}{t})U_4(t)Q_B^{(2)}U_4^1(t)U_2(\stackrel{~}{t})Q(F_t)+C(G_t)`$ (4.13)
where $`\stackrel{~}{t}=(3t1)/(3t)`$ as defined in eq. (4.12), and the function $`F_t(w)`$ is given as
$`F_t(w)={\displaystyle \frac{1}{4}}(f^2+f^2)^2,f(w;t)f_4(f_2(w;\stackrel{~}{t});t).`$ (4.14)
Note here that $`F_t(w)`$ is slightly different from $`F_2^{(l=2)}(w,t)`$ since the former now includes the contribution from $`U_2(\frac{l1}{l+1})|_{l=2}`$ in eq. (4.10). Let us write the function $`F_t(w)`$ explicitly,
$`F_t(w)={\displaystyle \frac{1}{4}}{\displaystyle \frac{\left(4\stackrel{~}{t}+(1+\stackrel{~}{t}^2)(w^2+w^2)\right)^2}{\left({\displaystyle \frac{9t^214t+9}{(t3)^2}}+2\stackrel{~}{t}w^2+\left({\displaystyle \frac{t+1}{t3}}\right)^2w^4\right)\left(ww^1\right)}}.`$ (4.15)
In the limit of $`t1`$, $`F_t(w)`$ becomes
$`F_t(w){\displaystyle \frac{1}{4}}(ww^1)^2,`$ (4.16)
in other words,
$`Q_B^{(2)}(t)Q_B^{(1)},`$
the result of the previous subsection.
Now how about the inverse? The relation $`U^1(t)Q_B^{(1)}(t)U(t)=Q_B^{(2)}`$ must hold since algebraically $`U^{(1)}(t)U(t)=1`$. Indeed, we will confirm this relation shortly. However, the relation is not the direct answer to our question. To find an answer, we would rather like to consider the operator $`U^1(t)Q_B^{(1)}(s)U(t)`$. If we could take the limit of $`s1`$ first, then take the other limit, the operator $`\underset{t1}{lim}U(t)`$ would be $`๐ฐ`$ we are looking for.
In calculating $`U^1(t)Q_B^{(1)}(s)U(t)`$, we first consider $`U_2(\stackrel{~}{t})Q_B^{(1)}(s)U_2^1(\stackrel{~}{t})`$. The function $`F_s(w)`$ in the integrand $`Q(F_s)`$ is replaced as<sup>5</sup><sup>5</sup>5The other function $`G_s(w)`$ is also changed accordingly.
$`F_s(w)F_s(z)|_{z=f_2(w;\stackrel{~}{t})}.`$ (4.17)
An explicit calculation of the rhs of eq. (4.17) shows that the factors in the numerator and denominator of the resultant expression carries terms with $`w^2`$ and $`w^4`$ other than constants. In the next step, we let the operator $`U_4(t)`$ act on the rhs of eq. (4.17) and the variable $`w`$ is replaced by $`f_4(w;t)`$, which has the fourth order branch cuts. The factor $`w^2`$ is replaced by $`\left(f_4(w;t)\right)^2`$ and it still has branch cuts. Generically those produce the ambiguity in defining the contour integration. In order to avoid this difficulty, we require the terms of $`w^2`$ vanish on the rhs of eq. (4.17): this condition is found to be $`s=t`$. If we take $`s=t`$, there is no ambiguity and we find that $`U^1(t)Q_B^{(1)}(t)U(t)`$ is certainly $`Q_B^{(2)}`$, even before taking the limit of $`t1`$.
We have not been able to make sense of taking the limits of $`U^1(t)Q_B^{(1)}(s)U(t)`$ independently. So we still have not reached the complete understanding of the relations between various TTK solutions.
### 4.3 Operator $`U(t)`$ and string field transformation
Though with some limitation, we have constructed the operator that relates the $`l`$-th and the first solutions of TTK solutions in the limit of $`t1`$. Here we describe some properties of the operator: the differential equation; some evidence that suggests the well-definedness of the operator. For concreteness, we consider the operator $`U(t)U_2^1(\stackrel{~}{t})U_4(t)`$ relating the first and the second solutions. With the operator $`U(t=1)`$, we are supposed to be able to relate the string fields defined around two solutions: $`|\mathrm{\Phi }^{(2)}=U(1)^1|\mathrm{\Phi }^{(1)}`$. Later, we will show the relation in terms of components fields.
In deriving the differential equation for $`U(t)`$, the parameter $`t`$ is to be restricted for $`|t|<1`$ as we will see shortly. It is straightforward to obtain the equation
$`{\displaystyle \frac{d}{dt}}U(t)={\displaystyle \frac{1}{2}}\left(K_2+{\displaystyle \frac{1}{2}}U_2(\stackrel{~}{t})K_4U_2^1(\stackrel{~}{t})\right){\displaystyle \frac{U(t)}{1t^2}}.`$ (4.18)
The operator $`U_2(\stackrel{~}{t})K_4U_2^1(\stackrel{~}{t})`$ may further be calculated as
$`U_2(\stackrel{~}{t})K_4U_2^1(\stackrel{~}{t})={\displaystyle \frac{dz}{2\pi i}z(z^4z^4)\left(\frac{df_2(z;\stackrel{~}{t})}{dz}\right)^2T\left(f_2(z;\stackrel{~}{t})\right)}.`$ (4.19)
where the integration path is over the unit circle. By changing the variable from $`z`$ to $`uf_2(z;\stackrel{~}{t})`$, we obtain the expression of the operator
$`U_2(\stackrel{~}{t})K_4U_2^1(\stackrel{~}{t})={\displaystyle \frac{du}{2\pi i}u\frac{(1+\stackrel{~}{t}^2)(u^4u^4)+4\stackrel{~}{t}(u^2u^2)}{(1+\stackrel{~}{t}^2u^2)(1+\stackrel{~}{t}^2u^2)}T(u)}.`$ (4.20)
It is easy to understand the integration path for $`u`$ is again over the unit circle, when the parameter $`t`$ is real and $`|t|<1`$. When we take $`t=1`$, the integration path is not clearly defined. The value is to be reached only in the limiting procedure. The integrand of (4.20) may be expanded with respect to $`\stackrel{~}{t}`$ since $`|\stackrel{~}{t}|<1`$ for $`|t|<1`$. We finally obtain the differential equation for $`U(t)`$,
$`{\displaystyle \frac{d}{dt}}U(t)`$ $`=`$ $`K(t)U(t),`$
$`K(t)`$ $``$ $`{\displaystyle \frac{(75t)(1+t)}{(1t)(3t)^3}}K_2+{\displaystyle \frac{16(1t^2)}{(t3)^4}}{\displaystyle \underset{2}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{13t}{3t}}\right)^{n2}K_{2n}.`$ (4.21)
Now, we show that $`U(t)^1`$ has a well-defined normal ordered expression:
$`U(t)^1=\mathrm{exp}\left({\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}v_{2n}L_{2n}\right)\mathrm{exp}\left(v_0L_0\right)\mathrm{exp}\left({\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}v_{2n}L_{2n}\right)`$ (4.22)
where the coefficients $`v_n=v_n(t)`$ are finite even in the limit of $`t1`$. The normal ordered form in eq. (4.22) are expressed with even modes Virasoro operators since $`U_2(t)`$ and $`U_{2l}(t)`$ themselves are written with $`L_{\pm 2}`$ and $`L_{\pm 2l}`$ and the algebra of even mode operators is closed.
In the following, we take the operator $`U(t)`$ relating the solutions $`l=1`$ and 2 as an example and confirm our claim. We have not noticed any difficulty to extend our analysis to other cases.
Taking $`|h`$, a highest weight state with the dimension $`h`$, we calculate various matrix elements of $`U(t)^1`$ in terms of $`v_n(t)`$ :
$`h\left|U(t)^1\right|h`$ $`=`$ $`e^{v_0h},`$ (4.23)
$`h\left|U(t)^1L_2\right|h`$ $`=`$ $`4hv_2e^{v_0h},`$ (4.24)
$`h\left|L_2U(t)^1\right|h`$ $`=`$ $`4hv_2e^{v_0h},`$ (4.25)
$`h\left|L_4U(t)^1\right|h`$ $`=`$ $`\left\{12h\left(v_2\right)^2+8hv_4\right\}e^{v_0h},`$ (4.26)
$`h\left|U(t)^1(L_2)^2\right|h`$ $`=`$ $`\left\{16h\left(v_2\right)^2+24hv_4\right\}e^{v_0h},`$ (4.27)
and so on. Then, if all of these matrix elements are obtained, we can determine the coefficients $`v_n(t)`$ iteratively. In appendix C, we calculate some of the matrix elements using the definition of $`U(t)^1`$. As a result, the coefficients $`v_n(t)`$ are found to be
$`v_0(t)`$ $`=`$ $`{\displaystyle \frac{1}{4}}\mathrm{ln}\left\{{\displaystyle \frac{64(1t)^3(3t)^2}{(9t^214t+9)^3}}\right\},`$ (4.28)
$`v_2(t)`$ $`=`$ $`{\displaystyle \frac{(6t^2+4t18)(3t1)}{4(9t^214t+9)(3t)}},`$ (4.29)
$`v_2(t)`$ $`=`$ $`{\displaystyle \frac{3t1}{2}}\left({\displaystyle \frac{1t}{9t^214t+9}}\right)^{\frac{1}{2}},`$ (4.30)
$`v_4(t)`$ $`=`$ $`{\displaystyle \frac{8t(1t)^2(9t^34t^211t+18)}{(t3)^2(9t^214t+9)^2}},`$ (4.31)
$`v_4(t)`$ $`=`$ $`{\displaystyle \frac{t}{4}}.`$ (4.32)
We should emphasize that there is no singularity in $`v_n(t)`$ in the limit $`t1`$. Finally, taking the limit, we can obtain the normal ordered expression of $`U(1)^1`$:
$`U(1)^1=\underset{t1}{lim}U(t)^1`$ (4.33)
$`=`$ $`\mathrm{exp}\left({\displaystyle \frac{1}{2}}L_2{\displaystyle \frac{1}{4}}L_4+\mathrm{}\right)\mathrm{exp}\left({\displaystyle \frac{1}{2}}\mathrm{log}2L_0\right)\mathrm{exp}\left({\displaystyle \frac{1}{8}}L_2+{\displaystyle \frac{1}{32}}L_4+\mathrm{}\right).`$
Now, let us consider the string field transformation $`|\mathrm{\Phi }^{(2)}=U(1)^1|\mathrm{\Phi }^{(1)}`$, by which two theories expanded around $`l=1`$ and $`l=2`$ solutions can be related. Since the operator $`U(1)^1`$ has the normal ordered expression (4.33), the string field transformation has a well-defined Fock space expression, namely we can obtain transformations for all component fields without any divergence.
Write the string field up to level two as
$`|\mathrm{\Phi }^{(1)}`$ $`=`$ $`\varphi (x)c_1|0`$ (4.34)
$`+A_\mu (x)c_1\alpha _1^\mu |0+iB(x)c_0|0`$
$`+\psi _{\mu \nu }(x)c_1\alpha _1^\mu \alpha _1^\nu |0+ia_\mu (x)c_1\alpha _2^\mu |0`$
$`+s(x)c_1|0+t(x)c_0c_1b_2|0+iu_\mu (x)c_0\alpha _1^\mu |0+\mathrm{}.`$
Acting the normal ordered expression (4.33) of $`U(1)^1`$ on the string field (4.34),<sup>6</sup><sup>6</sup>6We have used the commutation relations $`[L_m,\alpha _n^\mu ]=n\alpha _{m+n}^\mu `$, $`[L_m,c_n]=(2m+n)c_{m+n}`$, $`[L_m,b_n]=(mn)b_{m+n}`$, and $`[L_m,\phi (x)]=i\sqrt{2\alpha ^{}}_\mu \phi (x)\alpha _m^\mu (m0)`$ for any component fields $`\phi (x)`$. we can easily find transformations for these component fields:
$`\varphi ^{}(x)`$ $`=`$ $`\sqrt{2}e^{\frac{\alpha ^{}}{2}\mathrm{log}2^2}\left(\varphi (x)+{\displaystyle \frac{1}{8}}\psi _\mu ^\mu (x)+{\displaystyle \frac{\sqrt{2\alpha ^{}}}{4}}_\mu a^\mu (x){\displaystyle \frac{3}{8}}s(x)+{\displaystyle \frac{1}{2}}t(x)\mathrm{}\right),`$
$`A_\mu ^{}(x)`$ $`=`$ $`e^{\frac{\alpha ^{}}{2}\mathrm{log}2^2}\left(A_\mu (x)+\mathrm{}\right),`$
$`B^{}(x)`$ $`=`$ $`e^{\frac{\alpha ^{}}{2}\mathrm{log}2^2}\left(B(x)+\mathrm{}\right),`$
$`\psi _{\mu \nu }^{}(x)`$ $`=`$ $`e^{\frac{\alpha ^{}}{2}\mathrm{log}2^2}\{{\displaystyle \frac{\sqrt{2}}{4}}g_{\mu \nu }\varphi (x)+{\displaystyle \frac{1}{\sqrt{2}}}\psi _{\mu \nu }(x)`$
$`{\displaystyle \frac{\sqrt{2}}{4}}g_{\mu \nu }({\displaystyle \frac{1}{8}}\psi _\rho ^\rho (x)+{\displaystyle \frac{\sqrt{2\alpha ^{}}}{4}}_\rho a^\rho (x){\displaystyle \frac{3}{8}}s(x)+{\displaystyle \frac{1}{2}}t(x))+\mathrm{}\},`$
$`a_\mu ^{}(x)`$ $`=`$ $`e^{\frac{\alpha ^{}}{2}\mathrm{log}2^2}\{\sqrt{\alpha ^{}}_\mu \varphi (x)+{\displaystyle \frac{1}{\sqrt{2}}}a_\mu (x)`$
$`+\sqrt{\alpha ^{}}_\mu ({\displaystyle \frac{1}{8}}\psi _\rho ^\rho (x)+{\displaystyle \frac{\sqrt{2\alpha ^{}}}{4}}_\rho a^\rho (x){\displaystyle \frac{3}{8}}s(x)+{\displaystyle \frac{1}{2}}t(x))+\mathrm{}\},`$
$`s^{}(x)`$ $`=`$ $`e^{\frac{\alpha ^{}}{2}\mathrm{log}2^2}\{{\displaystyle \frac{3\sqrt{2}}{2}}\varphi (x)+{\displaystyle \frac{1}{\sqrt{2}}}s(x)`$
$`{\displaystyle \frac{3\sqrt{2}}{2}}({\displaystyle \frac{1}{8}}\psi _\rho ^\rho (x)+{\displaystyle \frac{\sqrt{2\alpha ^{}}}{4}}_\rho a^\rho (x){\displaystyle \frac{3}{8}}s(x)+{\displaystyle \frac{1}{2}}t(x))+\mathrm{}\},`$
$`t^{}(x)`$ $`=`$ $`e^{\frac{\alpha ^{}}{2}\mathrm{log}2^2}\{\sqrt{2}\varphi (x)+{\displaystyle \frac{1}{\sqrt{2}}}t(x)`$
$`+\sqrt{2}({\displaystyle \frac{1}{8}}\psi _\rho ^\rho (x)+{\displaystyle \frac{\sqrt{2\alpha ^{}}}{4}}_\rho a^\rho (x){\displaystyle \frac{3}{8}}s(x)+{\displaystyle \frac{1}{2}}t(x))+\mathrm{}\},`$
$`u_\mu ^{}(x)`$ $`=`$ $`e^{\frac{\alpha ^{}}{2}\mathrm{log}2^2}\left({\displaystyle \frac{1}{\sqrt{2}}}u_\mu (x)+\mathrm{}\right),`$
where the abbreviation denotes contributions from the higher level component fields.
The string field transformation has a well-defined expression. It mixes tensor fields of various ranks; On each component, it is a non-local transformation due to infinite derivative terms.
## 5 Discussion
In this paper, we addressed the question how the presence of many analytic classical solutions for SFT could be consistent with the physical picture of the tachyon condensation. Our result suggests that they are related by a particular type of gauge transformations. In more concrete terms, we have seen that the BRS charge for the $`l`$-th classical solution $`(l1)`$ can be transformed down to that for $`l=1`$. The inverse operation has some subtlety as explained in section 4. The transformation is generated by operators $`K_n(n=\mathrm{even})`$. We observed that the symmetry generated by operators $`K_n`$ are to be regarded as the โglobalโ part of the SFT gauge symmetry. The situation is summarized in the following sequence,
$`\{\mathrm{The}\mathrm{stringy}\mathrm{gauge}\mathrm{symmetry}:\mathrm{\Psi }^{}=U^1Q_\mathrm{B}U+U^1\mathrm{\Psi }U\}`$
$`\{\mathrm{Its}\mathrm{global}\mathrm{subset}:\mathrm{\Psi }^{}=U^1\mathrm{\Psi }U\mathrm{with}Q_\mathrm{B}U=0\}`$
$`\{\mathrm{The}\mathrm{symmetry}\mathrm{generated}\mathrm{with}K_n:\mathrm{\Psi }^{}=\mathrm{exp}(K(v))\mathrm{\Psi }\}`$
$`\{\mathrm{The}\mathrm{symmetry}\mathrm{generated}\mathrm{with}K_n(n=\mathrm{even})\}.`$
In relating TTK solutions, we have utilized the last subset in the above sequence. Generically speaking, solutions are to be related by the gauge symmetry. So our approach may be too restrictive and that could be the reason why we encounter the subtlety.
In order to confirm that the operator relating BRS charges is well-defined, we studied properties of the operator that transforms $`l=2`$ BRS charge into $`l=1`$ charge and found that it has a well-defined normal ordered expression in terms of the Virasoro generators.
We studied relations between solutions obtained in Refs. . Another important direction of investigation is to find how those solutions could be related to solutions obtained in different approaches, eg, the level truncation .
Most of the works on classical solutions for CSFT have been performed in the Siegel gauge. However, the universal solution proposed by cannot be in the Siegel gauge as explained in the appendix D. A transformation generated by $`K_n`$ cannot bring a universal solution into the Siegel gauge: we have to consider more general gauge transformation.
In relation to the VSFT conjecture and the TTK solutions, recently there appeared an interesting paper . The VSFT conjecture on the tachyon vacuum implies that the action expanded around a TTK solution must be related to VSFT via an appropriate transformation of the string field. The authors of discussed this possibility and constructed, with the level truncation technique and a regulated butterfly state, a classical solution that could clarify this point.
Before closing, let us add a few remarks. 1) The cohomology analysis around a classical solution has shown that the ghost numbers of non-trivial states depend on the value of $`l`$ . It would be interesting to see how operators, eg, $`U(t)`$ in section 4, relate cohomologically non-trivial states obtained for various $`l`$. That would be another non-trivial test of those operators and the question certainly deserves further study. 2) We wonder what happens to the symmetries of the SFT defined around the non-trivial classical solution. Here we make an observation that symmetries generated by $`K_n`$ are broken on these classical solutions. When we consider an infinitesimal change of the $`l`$-th solution in Ref. , the function $`F_{2l}(w)`$ is transformed as shown in eq. (3.16). It is easy to see that any choice of the parameters $`\epsilon _n`$ do not leave the function invariant. This implies the breaking of the symmetries: the symmetry generated by $`K_n`$ does not survive the tachyon condensation, at least, in its original form. 3) The function $`F(z)`$ in the generic form of classical solution (3.13) is to satisfy two conditions $`F(1/w)=F(w)`$ and $`F(\pm i)=0`$ and is related to another function $`G(z)`$ as in eq. (3.13). Strictly speaking, in order for (3.13) to be a classical solution, these conditions are required to hold only on the unit circle. We encounter the same situation in eq. (B.8).
## Acknowledgments
The authors are grateful to I. Kishimoto, T. Kugo and T. Huruya for discussions. This work is completed during our stay at the Yukawa Institute for Theoretical Physics at Kyoto University. We wish to thank for their kind hospitality extended to us. Discussions during the YITP workshop YITP-W-04-03 on โQuantum Field Theory 2004โ were useful to finish this work. This work is supported in part by the Grants-in-Aid for Scientific Research No. 13135209 and 15540262 from the Japan Society for the Promotion of Science.
## Appendix A Evaluating $`S[g^1Q_Bg,Q_B]`$ for $`g`$ connected to $`I`$
We consider the stringy gauge functional $`g`$ that may be continuously deformed to the identity $`I`$: ie, we assume that there exit a one-parameter family of functionals $`g(t)(0t1)`$ so that $`g(0)=I`$ and $`g(t=1)=g`$. Construct the pure gauge string field as
$`\mathrm{\Psi }(t)g^1(t)Q_Bg(t).`$ (A.1)
By using the properties (4), it is easy to show that the string field $`\mathrm{\Psi }(t)`$ satisfies the equation of motion (2.4). Note also that $`\mathrm{\Psi }(t=0)=0`$. Now we may calculate the variation of the action for $`\mathrm{\Psi }(t)`$ as
$`{\displaystyle \frac{d}{dt}}S[\mathrm{\Psi }(t),Q_B]={\displaystyle \frac{1}{g^2}}{\displaystyle \left(Q_B\mathrm{\Psi }(t)+\mathrm{\Psi }(t)\mathrm{\Psi }(t)\right)\frac{d}{dt}\mathrm{\Psi }(t)}.`$ (A.2)
The rhs of (A.2) is proportional to the equation of motion, and it vanishes. Obviously, it holds that $`S[\mathrm{\Psi }(t=0),Q_B]=0`$. Therefore $`S[\mathrm{\Psi }(t),Q_B]=0`$ for any value of $`t`$, in particular $`t=1`$ .
## Appendix B On the global transformation given in eq. (3.6)
The energy momentum tensor $`T(w)`$ is expanded by the Virasoro operator $`L_n`$ as
$`T(w)={\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}L_nw^{n1}.`$ (B.1)
From commutation relations of $`L_n`$, we can derive the commutation relation between $`T(w)`$ and $`T(w^{})`$ as
$`[T(w),T(w^{})]=T(w)\delta (w,w^{})+T(w)_w^{}\delta (w,w^{}),`$ (B.2)
where the delta function is defined by $`\delta (w,w^{})=_nw^nw_{}^{}{}_{}{}^{n1}`$. Here, we define half string operators associated with the energy-momentum tensor as follows,
$`๐ฏ_L(f)={\displaystyle _{C_{\mathrm{left}}}}{\displaystyle \frac{dw}{2\pi i}}f(w)T(w),๐ฏ_R(f)={\displaystyle _{C_{\mathrm{right}}}}{\displaystyle \frac{dw}{2\pi i}}f(w)T(w).`$ (B.3)
Using (B.2) and the splitting properties of the delta function , we can find the commutation relations between these operators:
$`[๐ฏ_L(f),๐ฏ_L(g)]`$ $`=`$ $`๐ฏ_L((f)gfg),`$ (B.4)
$`[๐ฏ_R(f),๐ฏ_R(g)]`$ $`=`$ $`๐ฏ_R((f)gfg),`$ (B.5)
$`[๐ฏ_L(f),๐ฏ_R(g)]`$ $`=`$ $`0,`$ (B.6)
where the functions $`f(w)`$ and $`g(w)`$ satisfy $`f(\pm i)=g(\pm i)=0`$.
If the function $`f(w)`$ satisfies $`f(w)=(dw/d\stackrel{~}{w})f(\stackrel{~}{w})`$ for $`\stackrel{~}{w}=1/w`$, we find that
$`dwf(w)T(w)=d\stackrel{~}{w}f(\stackrel{~}{w})\stackrel{~}{T}(\stackrel{~}{w})`$ (B.7)
since $`T(w)`$ is a primary field with the conformal dimension $`2`$ for $`c=0`$ . Using the relation (B.7), we can obtain two properties of the half string operators for $`f(w)`$ such that $`f(w)=(dw/d\stackrel{~}{w})f(\stackrel{~}{w})`$:
$`๐ฏ_R(f)AB=A๐ฏ_L(f)B,๐ฏ_R(f)I+๐ฏ_L(f)I=0,`$ (B.8)
where $`A`$ and $`B`$ denote arbitrary string fields and $`I`$ is the identity string field. In deriving eq. (B.8), it is suffice for the function $`f(w)`$ to satisfy the condition, $`f(w)=(dw/d\stackrel{~}{w})f(\stackrel{~}{w})`$ for $`\stackrel{~}{w}=1/w`$, on the unit circle.
We consider gauge transformation with the string functional
$`g=\mathrm{exp}\left(๐ฏ_L(f)I\right),`$ (B.9)
where $`f(w)`$ is required to satisfy the condition $`f(w)=(dw/d\stackrel{~}{w})f(\stackrel{~}{w})`$. It is easy to see that the other condition $`f(\pm i)=0`$ follows from the former. The gauge transformation is given by
$`\mathrm{\Psi }^{}=g^1Q_\mathrm{B}g+g^1\mathrm{\Psi }g.`$ (B.10)
The first term becomes zero since $`[Q_\mathrm{B},๐ฏ_L(f)]=0`$ and $`Q_\mathrm{B}I=0`$. Then, from the equations (B.8), the gauge transformation can be rewritten as the string field redefinition
$`\mathrm{\Psi }^{}=\mathrm{exp}(๐ฏ(f))\mathrm{\Psi },`$ (B.11)
where the operator $`๐ฏ(f)`$ is defined as $`๐ฏ(f)=๐ฏ_L(f)+๐ฏ_R(f)`$.
## Appendix C Matrix elements of $`U(t)^1`$
The operators $`U_2(\stackrel{~}{t})`$ and $`U_4(t)^1`$ can be written using normal ordered expression as
$`U_2(\stackrel{~}{t})=e^{\frac{\stackrel{~}{t}}{2}L_2}e^{\frac{1}{2}\mathrm{log}(1\stackrel{~}{t}^2)L_0}e^{\frac{\stackrel{~}{t}}{2}L_2},U_4(t)^1=e^{\frac{t}{4}L_4}e^{\frac{1}{4}\mathrm{log}(1t^2)L_0}e^{\frac{t}{4}L_4}.`$ (C.1)
Take a normalized highest weight state with the dimension $`h`$, $`|h`$, and calculate the matrix element $`h\left|U(t)^1\right|h`$,
$`h\left|U(t)^1\right|h=h\left|U_4(t)^1U_2(\stackrel{~}{t})\right|h=(1t^2)^{\frac{h}{4}}(1\stackrel{~}{t}^2)^{\frac{h}{2}}h\left|e^{\frac{t}{4}L_4}e^{\frac{\stackrel{~}{t}}{2}L_2}\right|h.`$ (C.2)
We can derive the recursion relation for matrix elements $`h\left|(L_4)^n(L_2)^{2n}\right|h`$,
$`h\left|(L_4)^n(L_2)^{2n}\right|h=8n(2n1)(4n+3h4)h\left|(L_4)^{n1}(L_2)^{2(n1)}\right|h,`$ (C.3)
which can be solved to give the expression
$`h\left|(L_4)^n(L_2)^{2n}\right|h=64^nn!{\displaystyle \frac{\mathrm{\Gamma }\left(n+\frac{1}{2}\right)}{\mathrm{\Gamma }\left(\frac{1}{2}\right)}}{\displaystyle \frac{\mathrm{\Gamma }\left(n+\frac{3h}{4}\right)}{\mathrm{\Gamma }\left(\frac{3h}{4}\right)}}.`$ (C.4)
Using eq. (C.4), we obtain
$`h\left|e^{\frac{t}{4}L_4}e^{\frac{\stackrel{~}{t}}{2}L_2}\right|h=(1+t\stackrel{~}{t}^2)^{\frac{3h}{4}}.`$ (C.5)
Substituting (C.5) into (C.2), we reach the final expression for $`h\left|U(1)^1\right|h`$:
$`h\left|U(t)\right|h=\left\{{\displaystyle \frac{64(1t)^3(3t)^2)}{(9t^214t+9)^3}}\right\}^{\frac{h}{4}}.`$ (C.6)
Next, let us calculate the matrix element $`h\left|U(t)^1L_2\right|h`$.
$`h\left|U(t)^1L_2\right|h`$ $`=`$ $`(1t^2)^{\frac{h}{4}}(1\stackrel{~}{t}^2)^{\frac{h+2}{2}}h\left|e^{\frac{t}{4}L_4}e^{\frac{\stackrel{~}{t}}{2}L_2}L_2\right|h`$ (C.7)
$`2h\stackrel{~}{t}(1t^2)^{\frac{h}{4}}(1\stackrel{~}{t}^2)^{\frac{h}{2}}h\left|e^{\frac{t}{4}L_4}e^{\frac{\stackrel{~}{t}}{2}L_2}\right|h.`$
Differentiating eq. (C.5) with respect to $`\stackrel{~}{t}`$, we find
$`h\left|e^{\frac{t}{4}L_4}e^{\frac{\stackrel{~}{t}}{2}L_2}L_2\right|h=3ht\stackrel{~}{t}(1+t\stackrel{~}{t}^2)^{\frac{3h}{4}1}.`$ (C.8)
Combining the results (C.5) and (C.8) with (C.7), we find
$`h\left|U(t)^1L_2\right|h`$ $`=`$ $`h{\displaystyle \frac{\stackrel{~}{t}(t\stackrel{~}{t}^23t2)}{1+t\stackrel{~}{t}^2}}\left\{{\displaystyle \frac{(1t^2)(1\stackrel{~}{t}^2)^2}{(1+t\stackrel{~}{t}^2)^3}}\right\}^{\frac{h}{4}}`$ (C.9)
$`=`$ $`h{\displaystyle \frac{2(3t^2+2t9)(3t1)}{(9t^214t+9)(3t)}}\left\{{\displaystyle \frac{64(1t)^3(3t)^2)}{(9t^214t+9)^3}}\right\}^{\frac{h}{4}}.`$
Other matrix elements may be calculated in a similar manner. Using the normal ordered expression of $`U_2(\stackrel{~}{t})`$ and $`U_4(t)`$, we easily find
$`h\left|L_2U(t)^1\right|h`$ $`=`$ $`(1t^2)^{\frac{h+2}{4}}(1\stackrel{~}{t}^2)^{\frac{h}{2}}`$ (C.10)
$`\times {\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(2n+1)!n!}}\left({\displaystyle \frac{\stackrel{~}{t}}{2}}\right)^{2n+1}({\displaystyle \frac{t}{4}})^nh\left|L_2(L_4)^n(L_2)^{2n+1}\right|h,`$
$`h\left|L_4U(t)^1\right|h`$ $`=`$ $`(1t^2)^{\frac{h+4}{4}}(1\stackrel{~}{t}^2)^{\frac{h}{2}}h\left|L_4e^{\frac{t}{4}L_4}e^{\frac{\stackrel{~}{t}}{2}L_2}\right|h`$ (C.11)
$`+(1t^2)^{\frac{h}{4}}(1\stackrel{~}{t}^2)^{\frac{h}{2}}\mathrm{\hspace{0.17em}2}hth\left|e^{\frac{t}{4}L_4}e^{\frac{\stackrel{~}{t}}{2}L_2}\right|h,`$
$`h\left|U(t)^1(L_2)^2\right|h`$ $`=`$ $`(1t^2)^{\frac{h}{4}}(1\stackrel{~}{t}^2)^{\frac{h+4}{2}}h\left|e^{\frac{t}{4}L_4}e^{\frac{\stackrel{~}{t}}{2}L_2}(L_2)^2\right|h`$ (C.12)
$`4(h+1)(1t^2)^{\frac{h}{4}}\stackrel{~}{t}(1\stackrel{~}{t}^2)^{\frac{h+2}{2}}h\left|e^{\frac{t}{4}L_4}e^{\frac{\stackrel{~}{t}}{2}L_2}L_2\right|h`$
$`+4h(h+1)(1t^2)^{\frac{h}{4}}\stackrel{~}{t}^2(1\stackrel{~}{t}^2)^{\frac{h}{2}}h\left|e^{\frac{t}{4}L_4}e^{\frac{\stackrel{~}{t}}{2}L_2}\right|h.`$
We can calculate (C.10) by using
$`h\left|L_2(L_4)^n(L_2)^{2n+1}\right|h=4^{3n+1}hn!{\displaystyle \frac{\mathrm{\Gamma }\left(n+\frac{3}{2}\right)}{\mathrm{\Gamma }\left(\frac{3}{2}\right)}}{\displaystyle \frac{\mathrm{\Gamma }\left(n+\frac{3h}{4}+\frac{1}{2}\right)}{\mathrm{\Gamma }\left(\frac{3h}{4}+\frac{1}{2}\right)}},`$ (C.13)
and eqs. (C.11) and (C.12) can be evaluated by using eqs. (C.5), (C.8) and
$`h\left|L_4e^{\frac{t}{4}L_4}e^{\frac{\stackrel{~}{t}}{2}L_2}\right|h`$ $`=`$ $`3h\stackrel{~}{t}^2(1+t\stackrel{~}{t}^2)^{\frac{3h}{4}1},`$
$`h\left|e^{\frac{t}{4}L_4}e^{\frac{\stackrel{~}{t}}{2}L_2}(L_2)^2\right|h`$ $`=`$ $`3ht(1+t\stackrel{~}{t}^2)^{\frac{3h}{4}2}(22t\stackrel{~}{t}^2+3ht\stackrel{~}{t}^2).`$
Finally, we obtain the expressions for the matrix elements,
$`h\left|L_2U(t)^1\right|h`$ $`=`$ $`2h\stackrel{~}{t}\left({\displaystyle \frac{1t^2}{1+t\stackrel{~}{t}^2}}\right)^{\frac{1}{2}}\left\{{\displaystyle \frac{(1t^2)(1\stackrel{~}{t}^2)^2}{(1+t\stackrel{~}{t}^2)^3}}\right\}^{\frac{h}{4}}`$
$`=`$ $`2h{\displaystyle \frac{(3t1)(1t)^{\frac{1}{2}}}{(9t^214t+9)^{\frac{1}{2}}}}\left\{{\displaystyle \frac{64(1t)^3(3t)^2)}{(9t^214t+9)^3}}\right\}^{\frac{h}{4}},`$
$`h\left|L_4U(t)^1\right|h`$ $`=`$ $`h{\displaystyle \frac{2t+3\stackrel{~}{t}^2t^2\stackrel{~}{t}^2}{1+t\stackrel{~}{t}^2}}\left\{{\displaystyle \frac{(1t^2)(1\stackrel{~}{t}^2)^2}{(1+t\stackrel{~}{t}^2)^3}}\right\}^{\frac{h}{4}}`$
$`=`$ $`h{\displaystyle \frac{9t^3+17t^23t+3}{9t^214t+9}}\left\{{\displaystyle \frac{64(1t)^3(3t)^2)}{(9t^214t+9)^3}}\right\}^{\frac{h}{4}},`$
$`h\left|U(t)^1(L_2)^2\right|h`$ $`=`$ $`[16h(h+1)\left\{{\displaystyle \frac{\stackrel{~}{t}(t\stackrel{~}{t}^23t2)}{4(1+t\stackrel{~}{t}^2)}}\right\}^2`$ (C.14)
$`24h{\displaystyle \frac{t(1\stackrel{~}{t}^2)(2+t\stackrel{~}{t}^2)}{8(1+t\stackrel{~}{t}^2)^2}}\left]\right\{{\displaystyle \frac{(1t^2)(1\stackrel{~}{t}^2)^2}{(1+t\stackrel{~}{t}^2)^3}}\}^{\frac{h}{4}}`$
$`=`$ $`[16h(h+1)\left\{{\displaystyle \frac{(3t1)(3t^2+2t9)}{2(3t)(9t^214t+9)}}\right\}^2`$
$`24h{\displaystyle \frac{8t(1t)^2(9t^34t^211t+18)}{(t3)^2(9t^214t+9)^2}}\left]\right\{{\displaystyle \frac{64(1t)^3(3t)^2)}{(9t^214t+9)^3}}\}^{\frac{h}{4}}.`$
## Appendix D The universal solutions are not in the Siegel gauge
In this appendix, we show that the universal solution given in Ref. does not satisfy the Siegel gauge condition.
The universal solution is written in the following form
$`|\mathrm{\Psi }_0`$ $`=`$ $`Q_\mathrm{L}(F)|I+C_\mathrm{L}(G)|I,`$ (D.2)
$`G={\displaystyle \frac{(F)^2}{1+F}}.`$
Let us search for the function $`F(w)`$ which satisfies the Siegel gauge condition
$`0=b_0|\mathrm{\Psi }_0={\displaystyle _{\mathrm{C}_{\mathrm{left}}}}{\displaystyle \frac{dw}{2\pi i}}F(w)b_0J_\mathrm{B}(w)|I+{\displaystyle _{\mathrm{C}_{\mathrm{left}}}}{\displaystyle \frac{dw}{2\pi i}}G(w)b_0c(w)|I,`$ (D.3)
relying on the conformal technique. First note that the identity state $`|I`$ can be written as $`|I=U_{IhI}^1|0`$. That is, the state can be expressed with the operator for the conformal transformation,
$`IhI={\displaystyle \frac{w^21}{2w}}g(w).`$
Using $`Q_n|0=0(n0)`$ and
$`U_gJ_\mathrm{B}(w)U_g^1=[g(w)]^{+1}J_\mathrm{B}(g(w))=g(w){\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}(g(w))^{n1}Q_n,`$
we find
$`J_\mathrm{B}(w)|I=U_g^1{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}g(w)(g(w))^{n1}Q_n|0.`$
So we obtain the expression for the first term of eq. (D.2)
$`Q_\mathrm{L}(F)|I`$ $`=`$ $`U_g^1{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\alpha _nQ_n|0`$ (D.4)
where $`\alpha _n`$ is given as
$`\alpha _n={\displaystyle _{\mathrm{C}_{\mathrm{left}}}}{\displaystyle \frac{dw}{2\pi i}}F(w)(g(w))^{n1}g(w).`$ (D.5)
We rewrite the second term $`C_\mathrm{L}(G)|I`$ in a similar manner: since
$`U_gc(w)U_g^1=[g(w)]^1c(g(w))=[g(w)]^1{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}(g(w))^{n+1}c_n,`$
we obtain
$`C_\mathrm{L}(G)|I=U_g^1{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\beta _nc_n|0,`$ (D.6)
with $`\beta _n`$ given as
$`\beta _n{\displaystyle _{\mathrm{C}_{\mathrm{left}}}}{\displaystyle \frac{dw}{2\pi i}}G(w)(g(w))^{n+1}(g(w))^1.`$ (D.7)
From eqs. (D.4) and (D.6), the universal solution is now rewritten as
$`|\mathrm{\Psi }_0=U_g^1\left\{{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\alpha _nQ_n|0+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\beta _nc_n|0\right\}.`$ (D.8)
Let us write the gauge condition for the universal solution written as eq. (D.8). First note
$`U_gb_0U_g^1={\displaystyle \frac{dw}{2\pi i}wU_gb(w)U_g^1}={\displaystyle \frac{dw}{2\pi i}w(g(w))^2b(g(w))}={\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}\gamma _nb_n,`$ (D.9)
where $`\gamma _n`$ is
$`\gamma _n={\displaystyle \frac{dw}{2\pi i}w(g(w))^2(g(w))^{n2}}={\displaystyle \frac{dw}{2\pi i}w\left(\frac{1+w^2}{2w^2}\right)^2\left(\frac{w^21}{2w}\right)^{n2}}.`$ (D.10)
After some calculation, we find that $`\gamma _n`$ vanish for $`n1`$ and $`n=`$odd. Therefore the condition may be written
$`0=b_0|\mathrm{\Psi }_0=U_g^1\left\{{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\gamma _{2m}\alpha _nb_{2m}Q_n|0+{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\gamma _{2m}\beta _nb_{2m}c_n|0\right\},`$
or,
$`{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\gamma _{2m}\alpha _nb_{2m}Q_n|0+{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\gamma _{2m}\beta _nb_{2m}c_n|0=0.`$ (D.11)
Explicitly writing the state $`b_{2m}Q_n|0(m1,n1)`$ as
$`b_{2m}Q_n|0=b_{2m}c_0L_n^X|0+\mathrm{},`$
we realize that the first and second terms in eq. (D.11) cannot cancel with each other. Thus $`\alpha _n=0`$ ($`n1`$) as well as $`\beta _n=0`$ ($`n1`$). Rewriting (D.5) with $`w=e^{i\sigma }`$, we find
$`\alpha _n=i^{n1}{\displaystyle _{\frac{\pi }{2}}^{\frac{\pi }{2}}}{\displaystyle \frac{d\sigma }{2\pi }}F(\sigma )\left(\mathrm{sin}\sigma \right)^{n1}\mathrm{cos}\sigma ={\displaystyle _1^1}{\displaystyle \frac{dx}{2\pi }}\stackrel{~}{F}(x)x^{n1}`$
In the last expression, we changed the variable as $`x=\mathrm{sin}\sigma `$ and used the notation $`F(\sigma )=\stackrel{~}{F}(x)`$. Clearly, the Siegel gauge condition requires the vanishing of the function, $`\stackrel{~}{F}(x)=0`$, therefore $`|\mathrm{\Psi }_0=0`$.
In conclusion, the universal functions cannot be the Siegel gauge.
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# 1 Introduction
## 1 Introduction
The analysis of (2+1) dimensional gravitational collapse models is useful as a toy model which may provide useful insights regarding the validity of Cosmic Censorship Hypothesis(CCH). There is no general proof for the validity of the hypothesis. There are many examples in which naked singularity is the outcome of the gravitational collapse. A good overview of the gravitational collapse scenario can be found in some of the review articles, Wald , Harada , Singh , Joshi . We can ask what the (2+1) dimensional Einsteinโs gravity has to say about CCH.
The study is important from the point of view of quantizing general relativity. The quantum gravitational effects play an important role in the gravitational collapse, especially when the collapse results in the formation of naked singularities. A lot of work has been done based on the canonical quantization of the LeMaitre-Tolman-Bondi(LTB) spherical dust collapse models . The canonical dynamics of the collapsing dust is analyzed by embedding the spherically symmetric ADM 4-metric in the LTB space time. Extending the quantization program to (2+1) dimensions is important due to the potential insights it might provide regarding the end state of the quantum gravitational collapse. This paper focuses mainly on obtaining the exact solutions of collapse of inhomogeneous spherical dust in (2+1) Einsteinโs theory of gravity, which can be useful for analyzing the quantum gravitational analogue.
The (2+1) dimensional gravity is fascinating in itโs own right. In this case there is no gravity outside matter. Matter curves space time only locally. Consequently there are no gravitational waves. The correspondence of Einsteinโs theory with Newtonian gravity breaks down. Newtonian gravity cannot be obtained as a limit of Einsteinโs theory. In the absence of a cosmological constant a dust distribution moves without any geodesic deviation between particles . These features make the (2+1) dimensional gravity an interesting toy model.
The work on homogeneous collapse of dust in (2+1) dimensions with a cosmological constant was done by Ross and Mann . They have shown that the stationary black hole solution found in arises naturally from the gravitational collapse of pressureless dust for a negative cosmological constant. Itโs properties were shown to be similar to the higher dimensional Oppenheimer-Snyder case. They have shown that there is no black hole formation for the case when cosmological constant is zero, agreeing with the results obtained in . They have also shown that the collapse to a naked singularity is possible for the positive cosmological constant case provided that the initial density is sufficiently small. We extend the analysis for inhomogeneous spherically symmetric dust distribution starting from generic initial velocity distribution. We show that the introduction of inhomogeneities in the initial density profile do not alter the qualitative features of the homogeneous case. This is drastically different from the (3+1) dimensional case in which the introduction of inhomogeneities in the initial density profile alter the nature of singularity .
We assume spherically symmetric dust. For the case of dust we set the pressure to zero. So the energy momentum tensor for dust is given by
$$T_{\mu \nu }=\rho u_\mu u_\nu ,$$
(1)
where $`\rho `$ is the density of the dust. We set up comoving coordinate system in which the three velocity of dust is given by $`u^\mu =(1,0,0)`$. We assume the comoving spherically symmetric metric of the form
$$ds^2=dt^2+e^{2b(t,r)}dr^2+R(t,r)^2d\varphi ^2$$
(2)
where $`t`$ is the coordinate time, $`r`$ is the label of the comoving shell. $`R(t,r)`$ represents the physical radius of the collapsing shell. $`\varphi `$ is the angular coordinate. The Einstein equations are
$$G_{\mu \nu }+\lambda g_{\mu \nu }=\kappa T_{\mu \nu }$$
(3)
where positive and negative $`\lambda `$ corresponds to three dimensional de-Sitter space and anti-deSitter space respectively. The non zero components of Einstein tensor are
$$G_{tt}=\frac{e^{2b}(b^{}R^{}R^{\prime \prime }+e^{2b}\dot{b}\dot{R})}{R},$$
(4)
$$G_{tr}=\frac{R^{}\dot{b}\dot{R^{}}}{R},$$
(5)
$$G_{rr}=\frac{e^{2b}\ddot{R}}{R},$$
(6)
$$G_{\varphi \varphi }=R^2(\dot{b}^2+\ddot{b}).$$
(7)
where is partial derivative w.r.t. $`r`$ and $`\dot{}`$ is partial derivative w.r.t. $`t`$. We study the solutions for the three cases ($`\lambda `$ is zero, negative, positive) separately.
## 2 Case $`\lambda =0.`$
The Einstein equations for $`\lambda =0`$ are as follows.
$$\frac{e^{2b}(b^{}R^{}R^{\prime \prime }+e^{2b}\dot{b}\dot{R})}{R}=\kappa \rho ,$$
(8)
$$\frac{R^{}\dot{b}\dot{R^{}}}{R}=0,$$
(9)
$$\frac{e^{2b}\ddot{R}}{R}=0,$$
(10)
$$R^2(\dot{b}^2+\ddot{b})=0.$$
(11)
Solving (10) gives
$$R=c_1t+c_2.$$
(12)
where $`c_1`$ and $`c_2`$ are functions of $`r`$ alone. The functions should be at least $`C^2`$. The function $`c_2`$ can be taken to be $`r`$ making use of the scaling freedom. This implies when $`t=0,R=r`$. $`c_1`$ is interpreted as the initial velocity of the particular shell with label $`r`$. The counter intuitive nature of (2+1) dimensions is already apparent from the equation because if one sets the initial velocity to zero, the entire dust cloud remains at rest without collapsing. The dust cloud always moves with uniform velocity. The behavior of dust is as if gravity were absent. This curious property can be explained because of the absence of any geodesic deviation between dust particles in (2+1) dimensions. This is related to the fact that the Riemann tensor is completely determined by the local matter distribution . Solving equation (11) gives
$$e^b=k_1t+k_2.$$
(13)
where $`k_1`$ and $`k_2`$ are functions of $`r`$ alone. One has to express these two arbitrary functions of $`r`$ in terms of the initial density and the initial velocity $`c_1`$ . Substituting (12) and (13) in (9) gives
$$c_1^{}k_2=k_1.$$
(14)
The conservation equation $`T_{;\nu }^{0\nu }=0`$ gives
$$\rho =\frac{\psi (r)}{e^bR},$$
(15)
where $`\psi (r)`$ is a function of $`r`$. Substituting equations (12) , (13) in (8) we get
$$\kappa \psi (k_1t+k_2)^2=(c_1^{}t+1)(k_1^{}t+k_2^{})(k_1t+k_2)c_1^{\prime \prime }t+k_1c_1(k_1t+k_2)^2.$$
(16)
Now collecting the coefficients of equal powers of $`t`$, the equation (16) can be brought to the form
$$A(r)t^2+B(r)t+C(r)=0.$$
(17)
This implies that $`A=0`$, $`B=0`$, $`C=0`$. It can be shown that each of the preceding three equations are interrelated. Each can be derived from the other two equations using the equation (14). If we write down $`C`$ explicitly we get
$$(\kappa \psi k_1c_1)k_2^2=k_2^{}.$$
(18)
From equation (15) we get
$$\psi =k_2r\rho _i,$$
(19)
where the symbol $`\rho _i`$ is used throughout the paper to represent the initial density profile, i.e $`\rho (0,r)`$. The function $`\rho _i(r)`$ is positive (in order to satisfy the weak energy condition) and is at least $`C^0`$. Substituting in (18) and integrating we get
$$\frac{1}{k_2^2}=c_1^22\kappa _0^r\rho _i(s)s๐s+c,$$
(20)
where $`c`$ is a positive constant. $`c`$ sets a natural scale of mass to the (2+1) dimensional universe. The constant $`\kappa `$ has the dimension of $`(mass)^1`$ . This feature is different from (3+1) dimensions where $`\kappa `$ can be made dimensionless . The metric (2) becomes
$$ds^2=dt^2+\frac{(c_1^{}t+1)^2dr^2}{c_1^22\kappa _0^r\rho _i(s)s๐s+c}+(c_1t+r)^2d\varphi ^2.$$
(21)
The functions $`c_1`$ and $`\rho _i`$ should be chosen such that $`c_1^22\kappa _0^r\rho _i(s)s๐s+c>0`$. A shell becomes singular only if initial velocity $`c_1`$ is negative. The time for singularity formation for a given shell is given by $`t=r/c_1`$, which is obtained by setting $`R`$, the physical radius of the shell to zero. The Ricci scalar $`R_i`$ diverges when the physical radius goes to zero.
$$R_i=\frac{2e^{2b}(b^{}R^{}R^{\prime \prime }+e^{2b}(\dot{b}\dot{R}+R(b^2+\ddot{b})+\ddot{R}))}{R}.$$
(22)
To answer questions about the nature of singularity obtained, we begin by looking for trapped surfaces. The shells are trapped if the divergence of a congruence of outgoing radial null geodesics $`\mathrm{\Theta }`$ is negative .
A congruence of outgoing radial null geodesics is considered having tangent vector $`(K^t,K^r,0)`$, where $`K^t=dt/dk`$ and $`K^r=dr/dk`$, where $`k`$ is an affine parameter along the geodesic which increases into the future. Imposing the conditions that $`K^t>0`$ and $`K^r>0`$, restricts the analysis to the case of future pointing outgoing null geodesic.
$$\mathrm{\Theta }=K_{;i}^i=\frac{1}{\sqrt{g}}\frac{}{x^i}(\sqrt{g}K^i).$$
(23)
For the metric (21) calculation of $`\mathrm{\Theta }`$ yields
$$\mathrm{\Theta }=\frac{K^r}{R}(e^b\dot{R}+R^{}).$$
(24)
Now we choose the initial velocity distribution such that there are no shell crossing singularities. This condition can be met if $`R^{}>0`$ at least till the time when singularity is formed for that particular $`r`$. This implies we have to choose the function $`c_1`$ such that $`c_1^{}/c_1<1/r`$. Now evaluating $`\mathrm{\Theta }`$ for the metric (21) yields
$$\mathrm{\Theta }=\frac{K^r(c_1^{}t+1)}{R}\left[\frac{c_1}{\sqrt{c_1^22\kappa _0^r\rho _i(s)s๐s+c}}+1\right].$$
(25)
Now in the above expression $`c_1^{}t+1=R^{}`$ which is greater than zero. We also have
$$\frac{K^t}{K^r}=e^b,$$
(26)
which is always positive if there are no shell crossing singularities. $`\mathrm{\Theta }`$ becomes negative only if the expression inside the square bracket becomes negative since $`K^r,R,R^{}`$ are positive. $`c_1`$ is negative since there would be no collapse if initial velocity is not negative. The condition for $`\mathrm{\Theta }`$ to become negative is
$$\frac{c_1}{\sqrt{c_1^22\kappa _0^r\rho _i(s)s๐s+c}}+1<0,$$
(27)
which implies
$$2\kappa _0^r\rho _i(s)s๐sc>0.$$
(28)
Now we call $`2_0^r\rho _i(s)s๐s`$ the mass function. It is a monotonically increasing function. The mass function for the shell $`r=0`$ is zero. Apparently, only those shells whose mass function is greater than $`c/\kappa `$ seem to be trapped. There is no dynamics as far as trapping is concerned since shells which are untrapped at time $`t=0`$ remain untrapped and vice versa. This rather unphysical property can be resolved when matching with an exterior is done.
The exterior solution for a circular symmetric (2+1) dimensional space time with a cosmological constant $`\lambda `$ is given by
$$ds^2=(\lambda R^2M)dT^2+\frac{dR^2}{(\lambda R^2M)}+R^2d\varphi ^2.$$
(29)
In the case when $`\lambda =0`$, $`M`$ should be negative in order to preserve the signature of the metric.
We set $`a=1/M`$ and rescale $`T`$ to $`\sqrt{1/a}T`$,
$$ds^2=dT^2+dR^2a+R^2d\varphi ^2.$$
(30)
We choose $`r_0`$ as the outer boundary of the collapsing dust. Matching involves equating the first and the second fundamental forms across $`r_0`$ . $`\varphi `$ is the same in the exterior and the interior metrics. Equating the coefficients of $`d\varphi ^2`$ on the hypersurface, we get
$$R_0=c_1(r_0)t+r_0.$$
(31)
Equating the remaining components
$$dt^2=dT^2+adR^2.$$
(32)
This implies
$$\frac{dT}{dt}=\sqrt{1+ac_1^2}.$$
(33)
The unit normal to the hypersurface in the exterior metric is $`(c_1\sqrt{a},\sqrt{a}\dot{T},0)`$. In the interior metric the unit normal to the hypersurface is $`(0,e^b,0)`$. Calculating the extrinsic curvature component $`K_{\varphi \varphi }^e`$ in the exterior metric yields
$$K_{\varphi \varphi }^e=R\sqrt{1/a+c_1^2}.$$
(34)
Calculating $`K_{\varphi \varphi }^i`$ in the interior gives
$$K_{\varphi \varphi }^i=(c_1t+r)\sqrt{c_1^22\kappa _0^{r_0}\rho _i(s)s๐s+c}.$$
(35)
Equating equations (34) and (35) we get
$$a=\frac{1}{2\kappa _0^{r_0}\rho _i(s)s๐s+c}.$$
(36)
This implies $`a`$ is positive only if $`c>2\kappa _0^{r_0}\rho _i(s)s๐s`$. So this means that $`c/\kappa `$ is the upper limit on the total mass function of the (2+1) dimensional collapsing star. Looking back at equation (28), the above restriction on the mass implies that all the shells are always untrapped.
We now analyze the nature of the singularity. The hypersurface $`R=0`$ (singularity curve) is given by $`t=r/c_1`$ in the $`tr`$ plane. The singularity curve is a monotonically increasing function (implied by the conditions that there are no shell crossing singularities). We identify a point on the singularity curve by $`P(t_s,r_s)`$. The point $`P`$ on the singularity curve is said to be at least locally naked if there exists an outgoing null geodesic which terminates in the past at $`P`$. We restrict attention to radial outgoing null geodesics. The equation of a radial null geodesic starting at a point $`P`$ on the singularity curve is obtained by imposing the null condition for the metric (21)
$$\frac{dt}{dr}=\pm \frac{(c_1^{}t+1)}{\sqrt{c_1^22\kappa _0^r\rho _i(s)s๐s+c}}.$$
(37)
The equation is of the form $`dt/dr=f(t,r)`$. Picardโs theorem states that if $`f(t,r)`$ and $`f(t,r)/t`$ are continuous functions on a closed rectangle $`G`$, then through each point $`(t_s,r_s)`$ in the interior of $`G`$ there passes a unique integral curve of the equation $`dt/dr=f(t,r)`$. The existence of a unique solution for the equation (37) is guaranteed by the Picardโs theorem if both $`f(t,r)`$ (given by R.H.S of (37)) and $`f/t`$ are continuous on the singularity curve and in the valid spacetime region of the $`tr`$ plane. For the continuity of $`f(t,r)`$ and $`f(t,r)/t`$, we require $`c_1`$, $`c_1^{}`$ and $`2\kappa _0^r\rho _i(s)s๐s`$ be continuous. The function $`c_1^{}`$ is at least $`C^1`$ since $`c_1`$ is at least $`C^2`$. The function $`\rho _i`$ is at least $`C^0`$ which implies $`2\kappa _0^r\rho _i(s)s๐s`$ is at least $`C^1`$. We also have $`c_1^22\kappa _0^r\rho _i(s)s๐s+c>0`$. This implies the continuity of $`f(t,r)`$ and $`f/t`$ in the valid spacetime as well as on the singularity curve. The null geodesic equation (37) can be readily solved. We get
$$t=\pm \left[e^{\pm {\scriptscriptstyle {\scriptscriptstyle \frac{c_1^{}dr}{\sqrt{c_1^22\kappa _0^r\rho _i(s)s๐s+c}}}}}\right]\left[\frac{e^{{\scriptscriptstyle {\scriptscriptstyle \frac{c_1^{}dr}{\sqrt{c_1^22\kappa _0^r\rho _i(s)s๐s+c}}}}}}{\sqrt{c_1^22\kappa _0^r\rho _i(s)s๐s+c}}+b\right],$$
(38)
where $`b`$ is a constant which is fixed by choosing the initial point. A radial null geodesic equation starting from the point $`P(t_s,r_s)`$ can be approximated by the Taylor expansion in the neighborhood of $`P`$
$$t=t_s+a_0(rr_s)+\frac{a_1(rr_s)^2}{2!}+\frac{a_2(rr_s)^3}{3!}+\mathrm{}..$$
(39)
The value of $`a_0`$ is equal to the R.H.S of (37) (with positive sign) evaluated at the point $`P(t_s,r_s)`$ ,
$$a_0=\frac{c_1^{}(r_s)t_s+1}{\sqrt{c_1^2(r_s)2\kappa _0^{r_s}\rho _i(s)s๐s+c}}.$$
(40)
Substituting $`t_s=r_s/c_1(r_s)`$, $`a_0`$ can be brought to the form
$$a_0=\frac{r_sc_1^{}(r_s)c_1(r_s)}{c_1^2(r_s)\left[\sqrt{1+\frac{c2\kappa _0^{r_s}\rho _i(s)s๐s}{c_1^2(r_s)}}\right]}.$$
(41)
Now we Taylor expand the singularity curve in the neighborhood of $`P(t_s,r_s)`$. We get
$$t=t_s+\frac{r_sc_1^{}(r_s)c_1(r_s)}{c_1^2(r_s)}(rr_s)+\mathrm{}$$
(42)
Now the singularity will be naked if $`a_0<(r_sc_1^{}(r_s)c_1(r_s))/c_1^2(r_s)`$. Physically this means that a null ray originating on the singularity curve at a shell with label $`r_s`$ is able to reach a neighboring shell $`r_s+dr`$ in the interval $`dt`$, before the shell $`r_s+dr`$ become singular. Now using the equation (42) we derive that $`a_0`$ is less than $`(r_sc_1^{}(r_s)c_1(r_s))/c_1^2(r_s)`$ if
$$c>2\kappa _0^{r_s}\rho _i(s)s๐s.$$
(43)
This implies that the outgoing null rays can emerge from the singularity till the critical shell (the shell for which $`c=2\kappa _0^r\rho _i(s)s๐s`$) becomes singular. From our earlier analysis regarding the validity of the exterior spacetime, we had obtained that the condition (43) should always be valid. This implies that the singularity is always naked. The hypersurface $`R=0`$ is timelike. From each point on the singularity curve we obtain a unique outgoing geodesic. This is different from the (3+1) dimensional case where a family of geodesics are shown to emerge in some cases from the central singularity ($`r=0`$) . This qualitative difference is due to the fact that Picardโs theorem is not applicable to the case considered in since $`f(t,r)/t`$ for the null geodesic equation $`dt/dr=f(t,r)`$ considered in is discontinuous at the central singularity. So the existence of family of solutions for the case considered in is not ruled out.
## 3 Case $`\lambda <0`$
Choose $`\mathrm{\Lambda }=\lambda `$. The Einstein equations give,
$$\ddot{R}+\mathrm{\Lambda }R=0,$$
(44)
$$\ddot{b}+\dot{b}^2+\mathrm{\Lambda }=0,$$
(45)
$$R^{}\dot{b}\dot{R}^{}=0,$$
(46)
$$\frac{e^{2b}(b^{}R^{}R^{\prime \prime }+e^{2b}(\dot{b}\dot{R}))}{R}+\mathrm{\Lambda }=\kappa \rho .$$
(47)
Solving equations (44) and (45) we get
$$R=A\mathrm{cos}(\sqrt{\mathrm{\Lambda }}t)+B\mathrm{sin}(\sqrt{\mathrm{\Lambda }}t),$$
(48)
$$e^b=D\mathrm{cos}(\sqrt{\mathrm{\Lambda }}t)+C\mathrm{sin}(\sqrt{\mathrm{\Lambda }}t),$$
(49)
where $`A,B,C,D`$ are functions of $`r`$ alone. The functions $`A,B`$ should be at least $`C^2`$. Exploiting the scaling freedom, $`A`$ can be set to $`r`$. So $`R=r`$ at time $`t=0`$. The function $`B`$ fixes the initial velocity profile. The functions $`D,C`$ have to be expressed in terms of the initial density and the initial velocity profile. Equation (46) implies $`R^{}/e^b`$ is a function of $`r`$. This implies
$$C=B^{}D.$$
(50)
By equation (15) we get $`\rho =\psi (r)/(re^b)`$. Simplification of equation (47) yields
$`(\mathrm{\Lambda }(Dr+CB)\kappa \psi )(D\mathrm{cos}(\sqrt{\mathrm{\Lambda }}t)+C\mathrm{sin}(\sqrt{\mathrm{\Lambda }}t))B^{\prime \prime }\mathrm{sin}(\sqrt{\mathrm{\Lambda }}t)`$
$$+\frac{(D^{}\mathrm{cos}(\sqrt{\mathrm{\Lambda }}t)+C^{}\mathrm{sin}(\sqrt{\mathrm{\Lambda }}t))(\mathrm{cos}(\sqrt{\mathrm{\Lambda }}t)+B^{}\mathrm{sin}(\sqrt{\mathrm{\Lambda }}t))}{D\mathrm{cos}(\sqrt{\mathrm{\Lambda }}t)+C\mathrm{sin}(\sqrt{\mathrm{\Lambda }}t)}=0.$$
(51)
This equation can be brought to the form
$$X(r)\mathrm{cos}^2(\sqrt{\mathrm{\Lambda }}t)+Y(r)\mathrm{cos}(\sqrt{\mathrm{\Lambda }}t)\mathrm{sin}(\sqrt{\mathrm{\Lambda }}t)+Z(r)\mathrm{sin}^2(\sqrt{\mathrm{\Lambda }}t)=0.$$
(52)
This equation is valid for all times,which implies $`X=0,Y=0,Z=0`$. Now $`X(r)`$ is
$$(\mathrm{\Lambda }(Dr+CB)\kappa \psi )D^2+D^{}=0.$$
(53)
Now $`\psi =\rho _irD`$ where $`\rho _i`$ is the initial density profile of the dust. Integrating equation (53) gives
$$\frac{1}{D^2}=\mathrm{\Lambda }r^2+\mathrm{\Lambda }B^22\kappa _0^r\rho _i(s)s๐s+c.$$
(54)
Constant $`c`$ has the same interpretation as the $`\mathrm{\Lambda }=0`$ case. The metric is
$$ds^2=dt^2+\frac{(\mathrm{cos}(\sqrt{\mathrm{\Lambda }}t)+B^{}\mathrm{sin}(\sqrt{\mathrm{\Lambda }}t))^2dr^2}{\mathrm{\Lambda }r^2+\mathrm{\Lambda }B^22\kappa _0^r\rho _i(s)s๐s+c}+(r\mathrm{cos}(\sqrt{\mathrm{\Lambda }}t)+B\mathrm{sin}(\sqrt{\mathrm{\Lambda }}t))^2d\varphi ^2.$$
(55)
The functions $`\rho _i`$ and $`B`$ should satisfy the inequality $`\mathrm{\Lambda }r^2+\mathrm{\Lambda }B^22\kappa _0^r\rho _i(s)s๐s+c>0`$. We choose the initial velocity profile such that there can be no shell crossing singularities at least till the time singularity is reached. The condition on $`B`$ is $`B^{}/B<1/r`$ (assuming that the initial velocity is negative). A shell with label $`r`$ becomes singular when the physical radius $`R`$ shrinks to zero size. The shell becomes singular at time $`t_s=\mathrm{arctan}(r/B)/\sqrt{\mathrm{\Lambda }}`$. To look for trapped surfaces we analyze $`\mathrm{\Theta }`$. Calculating $`\mathrm{\Theta }`$ using equation (24) for the metric (55) yields
$$\mathrm{\Theta }=\frac{K^r(\mathrm{cos}(\sqrt{\mathrm{\Lambda }}t)+B^{}\mathrm{sin}(\sqrt{\mathrm{\Lambda }}t))}{R}\left[1+\frac{\sqrt{\mathrm{\Lambda }}(r\mathrm{sin}(\sqrt{\mathrm{\Lambda }}t)+B\mathrm{cos}(\sqrt{\mathrm{\Lambda }}t))}{\sqrt{\mathrm{\Lambda }r^2+\mathrm{\Lambda }B^22\kappa _0^r\rho _i(s)s๐s+c}}\right].$$
(56)
For a shell with label $`r`$ to get trapped $`\mathrm{\Theta }<0`$, which implies the expression within square bracket becomes negative. The condition is
$$\mathrm{\Lambda }r^2\mathrm{cos}^2(\sqrt{\mathrm{\Lambda }}t)+\mathrm{\Lambda }B^2\mathrm{sin}^2(\sqrt{\mathrm{\Lambda }}t)+2\mathrm{\Lambda }rB\mathrm{cos}(\sqrt{\mathrm{\Lambda }}t)\mathrm{sin}(\sqrt{\mathrm{\Lambda }}t)2\kappa _0^r\rho _i(s)s๐s+c<0.$$
(57)
Now the sum of the first three terms on the L.H.S is equal to $`\mathrm{\Lambda }R^2`$. So the above inequality becomes
$$\frac{2\kappa _0^r\rho _i(s)s๐sc}{\mathrm{\Lambda }R^2}>1.$$
(58)
So this implies only those shells whose mass function is greater than $`c/\kappa `$ can ever become trapped. For the outer shells which can get trapped, trapping occurs when the the physical radius shrinks to a size $`R<\sqrt{(2\kappa _0^r\rho _i(s)s๐sc)/\mathrm{\Lambda }}`$. This is different from a (3+1) dimensional case where a black hole can be formed from any amount of mass by shrinking it to a radius less than the Scwharzschild radius. In (2+1) dimensional case if a star has the total mass function less than $`c/\kappa `$, there is no formation of trapped surfaces.
To analyze the nature of the singularity formed, we use the argument similar to the case $`\lambda =0`$. Now the hypersurface $`R=0`$ is the singularity curve in the $`tr`$ plane. It is given by the equation
$$t=\mathrm{arctan}(r/B)/\sqrt{\mathrm{\Lambda }}.$$
(59)
Now let $`P(t_s,r_s)`$ be a point on the singularity curve. The curve can be Taylor expanded near the point $`P`$. We get
$$t=t_s+\frac{r_sB^{}(r_s)B(r_s)}{\sqrt{\mathrm{\Lambda }}(B^2(r_s)+r_s^2)}(rr_s)+\mathrm{}.$$
(60)
The null condition for the metric (55) is
$$\frac{dt}{dr}=\frac{\mathrm{cos}(\sqrt{\mathrm{\Lambda }}t)+B^{}\mathrm{sin}(\sqrt{\mathrm{\Lambda }}t)}{\sqrt{\mathrm{\Lambda }r^2+\mathrm{\Lambda }B^22\kappa _0^r\rho _i(s)s๐s+c}}.$$
(61)
Picardโs theorem guarantees a unique solution for the geodesic starting from the point $`P(t_s,r_s)`$. Near the point $`P`$, we Taylor expand the solution,
$$t=t_s+a_0(rr_s)+\frac{a_1(rr_s)^2}{2!}+\mathrm{}.$$
(62)
The value of $`a_0`$ is the R.H.S of equation (61) evaluated at $`P(t_s,r_s)`$ where $`t_s=\mathrm{arctan}(r_s/B(r_s))/\sqrt{\mathrm{\Lambda }}`$. We obtain,
$$a_0=\frac{B^{}(r_s)r_sB(r_s)}{(B^2(r_s)+r_s^2)\sqrt{\mathrm{\Lambda }}\sqrt{\left[1+\frac{c2\kappa _0^{r_s}\rho _i(s)s๐s}{\mathrm{\Lambda }(B^2(r_s)+r_s^2)}\right]}}.$$
(63)
Comparing with equation (60) we obtain the condition that the singularity is naked if
$$a_0=\frac{B^{}(r_s)r_sB(r_s)}{(B^2(r_s)+r_s^2)\sqrt{\mathrm{\Lambda }}\sqrt{\left[1+\frac{c2\kappa _0^{r_s}\rho _i(s)s๐s}{\mathrm{\Lambda }(B^2(r_s)+r_s^2)}\right]}}<\frac{r_sB^{}(r_s)B(r_s)}{\sqrt{\mathrm{\Lambda }}(B^2(r_s)+r_s^2)}.$$
(64)
Now this is satisfied only if
$$c>2\kappa _0^{r_s}\rho _i(s)s๐s.$$
(65)
If we consider a star whose total mass function is greater than $`c/\kappa `$, we identify the critical shell $`r_c`$ for which $`_0^{r_c}\rho _i(s)s๐s=c/\kappa `$. Outgoing null rays can emerge from the singularity till the time when the critical shell becomes singular. This implies that during the collapse, the singularity will be at least locally naked till the time when the critical shell becomes singular. The $`R=0`$ hypersurface is therefore timelike for the shells $`r<r_c`$. It becomes null for $`r=r_c`$ and is spacelike for $`r>r_c`$.
The exterior in (2+1) dimensions is given by the BTZ metric with zero angular momentum,
$$ds^2=(\mathrm{\Lambda }R^2M)dT^2+\frac{dR^2}{(\mathrm{\Lambda }R^2M)}+R^2d\varphi ^2.$$
(66)
We denote the outer boundary of the dust by the coordinate $`r_0`$. Matching the metric components of (55) and (66) on the hypersurface $`r_0`$ gives
$$R_0=r_0\mathrm{cos}(\sqrt{\mathrm{\Lambda }}t)+B(r_0)\mathrm{sin}(\sqrt{\mathrm{\Lambda }}t),$$
(67)
$$1=(\mathrm{\Lambda }R^2M)\dot{T}^2\frac{\dot{R}^2}{(\mathrm{\Lambda }R^2M)}.$$
(68)
Unit normal to the hypersurface in the interior coordinates is $`(0,e^b,0)`$ and in the exterior coordinates is $`(\dot{R},\dot{T},0)`$. Calculating $`K_{\varphi \varphi }`$ in the interior gives
$$K_{\varphi \varphi }^i=e^bR(\mathrm{cos}(\sqrt{\mathrm{\Lambda }}t)+B^{}\mathrm{sin}(\sqrt{\mathrm{\Lambda }}t)).$$
(69)
Calculating $`K_{\varphi \varphi }`$ in the exterior gives
$$K_{\varphi \varphi }^e=\dot{T}R(\mathrm{\Lambda }R^2M).$$
(70)
Equating both we get
$$M=2\kappa _0^{r_0}\rho _i(s)s๐sc.$$
(71)
So in the exterior static solution, $`M`$ can be both positive and negative. The singularity will be naked if $`M`$ is negative otherwise it will be a black hole with event horizon at $`R=\sqrt{M/\mathrm{\Lambda }}`$.
## 4 Case $`\lambda >0`$
Solving the Einstein equations for the metric (2) with a positive $`\lambda `$ we get
$$R=A\mathrm{cosh}(\sqrt{\lambda }t)+B\mathrm{sinh}(\sqrt{\lambda }t),$$
(72)
$$e^b=D\mathrm{cosh}(\sqrt{\lambda }t)+C\mathrm{sinh}(\sqrt{\lambda }t),$$
(73)
where $`A,B,C,D`$ are functions of $`r`$ alone. $`A`$ can be set to $`r`$. So $`R=r`$ at time $`t=0`$. The function $`B`$ fixes the initial velocity profile. Solution of Einsteinโs equations proceeds much the same way as the negative $`\lambda `$ case. The metric therefore is
$$ds^2=dt^2+\frac{(\mathrm{cosh}(\sqrt{\lambda }t)+B^{}\mathrm{sinh}(\sqrt{\lambda }t))^2dr^2}{\lambda r^2+\lambda B^22\kappa _0^r\rho _i(s)s๐s+c}+(r\mathrm{cosh}(\sqrt{\lambda }t)+B\mathrm{sinh}(\sqrt{\lambda }t))^2d\varphi ^2.$$
(74)
There is a restriction on $`B`$ which comes from the fact that in the above metric, the signature has to be preserved.
$$\lambda r^2+\lambda B^22\kappa _0^r\rho _i(s)s๐s+c>0.$$
(75)
A shell with label $`r`$ becomes singular only if itโs initial velocity is sufficiently high to overcome the effect of positive cosmological constant. More precisely, when $`R`$ becomes zero , $`B=r/(\mathrm{tanh}\sqrt{\lambda }t_s)`$ where $`t_s`$ is the time for singularity formation. This implies that for singularity formation, it is necessary that $`B<r`$ for a given shell. If $`B>r`$ and $`B<0`$ then there will always be a rebounce at some finite $`t`$ and $`R`$ . The exterior is given by
$$ds^2=(\lambda R^2M)dT^2+\frac{dR^2}{(\lambda R^2M)}+R^2d\varphi ^2.$$
(76)
We get the right signature only if $`M`$ is negative. In this case the cosmological horizon occurs at $`R_c=\sqrt{M/\lambda }`$. We restrict our analysis to the case in which all shells collapse to singularity ($`B<r`$). We choose the outer shell $`r=r_0`$ such that the physical radius $`R`$ of the outer shell is less than the cosmological horizon $`R_c`$. The matching conditions on the hypersurface $`r=r_0`$ (outer boundary) is
$$M=2\kappa _0^{r_0}\rho _i(s)s๐sc.$$
(77)
Since $`M`$ is negative, $`c>2\kappa _0^{r_0}\rho _i(s)s๐s`$. The condition (75) gives $`B^2>r^2+M/\lambda `$. The condition $`B<r`$ which is required for the singularity formation, is more restrictive.
To analyze the nature of singularity, we calculate $`\mathrm{\Theta }`$ and evaluate the condition for trapping,
$$\mathrm{\Theta }=\frac{K^r(\mathrm{cosh}(\sqrt{\lambda }t)+B^{}\mathrm{sinh}(\sqrt{\lambda }t))}{R}[1+\frac{\sqrt{\lambda }(r\mathrm{sinh}(\sqrt{\lambda }t)+B\mathrm{cosh}(\sqrt{\lambda }t))}{\sqrt{\lambda r^2+\lambda B^22\kappa _0^r\rho _i(s)s๐s+c}}].$$
(78)
The condition for trapping is
$$\frac{c2\kappa _0^r\rho _i(s)s๐s}{\lambda R^2}<1.$$
(79)
Since $`c>2\kappa _0^r\rho _i(s)s๐s`$ the above condition for trapping is never met.
To see if there are outgoing null geodesics emerging from the singularity, we carry out the analysis as before. The singularity curve is given by $`t=\mathrm{tanh}^1(\mathrm{r}/\mathrm{B})/\sqrt{\lambda }`$. Expanding the curve in the neighborhood of a point $`P`$ on the singularity curve we get
$$t=t_s+\frac{r_sB^{}(r_s)B(r_s)}{\sqrt{\mathrm{\Lambda }}(B^2(r_s)r_s^2)}(rr_s)+\mathrm{}.$$
(80)
The null condition for the metric (74) is
$$\frac{dt}{dr}=\frac{\mathrm{cosh}(\sqrt{\lambda }t)+B^{}\mathrm{sinh}(\sqrt{\lambda }t)}{\sqrt{\lambda r^2+\lambda B^22\kappa _0^r\rho _i(s)s๐s+c}}.$$
(81)
Picardโs theorem guarantees a unique solution for the geodesic starting from the point $`P(t_s,r_s)`$. Near the point $`P`$, we approximate the solution by the series,
$$t=t_s+a_0(rr_s)+\frac{a_1(rr_s)^2}{2!}+\mathrm{}.$$
(82)
The value $`a_0`$ is the R.H.S of equation (81) evaluated at $`P(t_s,r_s)`$ where $`t_s=\mathrm{tanh}^1(r_s/B(r_s))/\sqrt{\lambda }`$. We obtain,
$$a_0=\frac{B^{}(r_s)r_sB(r_s)}{(B^2(r_s)r_s^2)\sqrt{\lambda }\sqrt{\left[1+\frac{c2\kappa _0^{r_s}\rho _i(s)s๐s}{\lambda (B^2(r_s)r_s^2)}\right]}}.$$
(83)
Comparing with equation (80) we obtain the condition that the singularity is naked if
$$a_0=\frac{B^{}(r_s)r_sB(r_s)}{(B^2(r_s)r_s^2)\sqrt{\lambda }\sqrt{\left[1+\frac{c2\kappa _0^{r_s}\rho _i(s)s๐s}{\lambda (B^2(r_s)r_s^2)}\right]}}<\frac{r_sB^{}(r_s)B(r_s)}{\sqrt{\lambda }(B^2(r_s)r_s^2)}.$$
(84)
For the above condition to be satisfied we require
$$\frac{c2\kappa _0^{r_s}\rho _i(s)s๐s}{\lambda (B^2(r_s)r_s^2)}>0.$$
(85)
For singularity formation, we derived earlier that $`B<r`$. So $`B^2>r^2`$. So the above inequality is met if
$$c>2\kappa _0^{r_s}\rho _i(s)s๐s.$$
(86)
Since this condition is always met in the collapse scenario under consideration, outgoing null rays can emerge from the singularity. The hypersurface $`R=0`$ is timelike. The singularity is therefore naked.
It is interesting to note that for $`3+1`$ dimensional collapse with positive $`\lambda `$, there are usually two horizons formed. One is the black hole event horizon and the other is the cosmological horizon .
## 5 Conclusions
Dust in (2+1) dimensions can undergo a variety of collapse scenarios. The curvature singularity formed due to the collapse of dust can be both naked or covered, depending on the sign of the cosmological constant and the initial density profile of the dust. In the absence of cosmological constant, collapse to a singularity is possible only if the initial velocity of dust is negative. The matching across the dust edge between the interior solution and the flat exterior is possible only if the total mass function of the collapsing dust ($`2_0^{r_0}\rho _i(s)s๐s`$, where $`r_0`$ is the outer boundary) has an upper bound. The analysis of $`\mathrm{\Theta }`$ (the divergence of outgoing null geodesics) rules out the formation of trapped surfaces during collapse. The singularity formed is naked since outgoing null rays are found to emerge from the singularity.
The collapse with a negative cosmological constant has several interesting features. In this scenario, collapse to a singularity is always possible. The analysis of $`\mathrm{\Theta }`$ indicates the existence of a lower bound on the mass function of the shells that can become trapped during the collapse. The exterior static solution (BTZ blackhole with zero angular momentum) can be matched to the interior metric. The matching across the dust edge gives the relationship between the parameter $`M`$ (analogous to the Schwarzschild mass) in the BTZ exterior, and the mass function of the outer shell. From the analysis of outgoing null geodesics emerging from the singularity, it is concluded that the singularity is either naked or covered depending on the mass function of the collapsing shell. The singularity is covered only if the mass function of the shell is greater than a critical value.
Collapse to a singularity occurs in the presence of a positive cosmological constant only if the initial velocity of the dust shell is high enough to overcome the repulsive effect of the positive cosmological constant. The collapse results in the formation of naked singularity.
## 6 Acknowledgment
I would like to thank T.P.Singh, Rakesh Tibrewala, Suresh Nampuri, Arun Madhav T Menon and Ashutosh Mahajan, for useful discussions.
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# Scattering Theory for Jacobi Operators with Quasi-Periodic Background
## 1. Introduction
Classical scattering theory deals with the reconstruction of a given Jacobi operator
(1.1)
$$Hu(n)=a(n)u(n+1)+a(n1)u(n1)+b(n)u(n),$$
which is a short range perturbation of the free one $`H_0`$ associated with the coefficients $`a(n)=\frac{1}{2}`$, $`b(n)=0`$. This case has been first developed on an informal level by Case in a series of papers . The first rigorous results were established by Guseinov , who gave necessary and sufficient conditions for the scattering data to determine $`H`$ uniquely under the assumption
(1.2)
$$\underset{n}{}|n|\left(|a(n)\frac{1}{2}|+|b(n)|\right)<\mathrm{}.$$
Further extensions were made by Guseinov , , and Teschl . Additional details and further references can be found, e.g., in .
In addition to being of interest on its own, scattering theory can also be used to solve the initial value problem for the Toda equation via the inverse scattering transform. This has been formally developed by Flaschka (see also and for the case of rapidly decaying sequences) who also worked out the inverse procedure in the reflection-less case. Further results and an extension of the method to the entire Toda hierarchy were given by Teschl in and .
The next interesting problem is to replace the free Hamiltonian $`H_0`$ by one with a periodic potential. First results in the case of Sturm-Liouville operators have been obtained by Firsova in a series of papers (see ). For further results, including potentials with different spatial asymptotics, and additional references see Gesztesy et al. . In the discrete case, the investigation has only recently been started by Boutet de Monvel and Egorova and by Volberg and Yuditskii , who treat the case where $`H`$ has a homogeneous spectrum and is of Szegรถ class exhaustively from an operator point of view. Applications to the Toda lattice can be found in Bazargan and Egorova and Boutet de Monvel and Egorova .
Finally, let us give a brief overview of the paper:
Section 2 collects some well-known facts from Riemann surfaces and introduces the necessary notation. Section 3 introduces the Baker-Akhiezer function and investigates the quasi-momentum map. In the periodic case, where the integrals can be explicitly computed, this was first done in . In addition, we characterize the second solution at the band edges. In Section 4 we prove existence of Jost solutions and use them to characterize the spectrum of the perturbed operator. In the periodic case, existence of Jost solutions was first shown by Geronimo and Van Assche and the fact that there are only finitely many eigenvalues in each gap was first proven in Cojuhari and later rediscovered in Teschl . Section 5 introduces the transformation operator and proves the crucial decay estimate on its coefficients. This was first done by Boutet de Monvel and Egorova in the periodic case under the additional assumption that all spectral gaps are open. We fix a problem in the original proof and at the same time simplify and streamline the argument. Section 6 investigates the scattering matrix. Our main result here is the reconstruction of the transmission coefficient from the reflection coefficient, which was not known previously, even in the periodic case. Section 7 derives the Gelโfand-Levitan-Marchenko equation and proves positivity of the Gelโfand-Levitan-Marchenko operator. In addition, we formulate necessary conditions for the scattering data to uniquely determine our Jacobi operator. Our final Section 8 shows that our necessary conditions for the scattering data are also sufficient. It should be mentioned that, due to the lack of continuity with respect to the spacial variable $`n`$, a significant change in the strategy of the original proof in the continuous case from is needed.
Our approach uses heavily the fact that the Baker-Akhiezer function is a meromorphic function on the Riemann surface associated with the problem. This strategy gives a more streamlined treatment and more elegant proofs even in the special cases which were previously known. In this respect it is important to emphasize that, in contradistinction to the constant background case, the upper sheet of our Riemann surface is not simply connected and in particular not isomorphic to the unit disc.
## 2. Quasi-periodic finite-gap operators and Riemann surfaces
To set the stage let $`๐`$ be the Riemann surface associated with the following function
(2.1)
$$R_{2g+2}^{1/2}(z),R_{2g+2}(z)=\underset{j=0}{\overset{2g+1}{}}(zE_j),E_0<E_1<\mathrm{}<E_{2g+1},$$
$`g`$. $`๐`$ is a compact, hyperelliptic Riemann surface of genus $`g`$. We will choose $`R_{2g+2}^{1/2}(z)`$ as the fixed branch
(2.2)
$$R_{2g+2}^{1/2}(z)=\underset{j=0}{\overset{2g+1}{}}\sqrt{zE_j},$$
where $`\sqrt{.}`$ is the standard root with branch cut along $`(\mathrm{},0)`$.
A point on $`๐`$ is denoted by $`p=(z,\pm R_{2g+2}^{1/2}(z))=(z,\pm )`$, $`z`$, or $`p=\mathrm{}_\pm `$, and the projection onto $`\{\mathrm{}\}`$ by $`\pi (p)=z`$. The points $`\{(E_j,0),0j2g+1\}๐`$ are called branch points and the sets
(2.3)
$$\mathrm{\Pi }_\pm =\{(z,\pm R_{2g+2}^{1/2}(z))z\backslash \underset{j=0}{\overset{g}{}}[E_{2j},E_{2j+1}]\}๐$$
are called upper, lower sheet, respectively.
Let $`\{a_j,b_j\}_{j=1}^g`$ be loops on the surface $`๐`$ representing the canonical generators of the fundamental group $`\pi _1(๐)`$. We require $`a_j`$ to surround the points $`E_{2j1}`$, $`E_{2j}`$ (thereby changing sheets twice) and $`b_j`$ to surround $`E_0`$, $`E_{2j1}`$ counter-clock wise on the upper sheet, with pairwise intersection indices given by
(2.4)
$$a_ia_j=b_ib_j=0,a_ib_j=\delta _{ij},1i,jg.$$
The corresponding canonical basis $`\{\zeta _j\}_{j=1}^g`$ for the space of holomorphic differentials can be constructed by
(2.5)
$$\underset{ยฏ}{\zeta }=\underset{j=1}{\overset{g}{}}\underset{ยฏ}{c}(j)\frac{\pi ^{j1}d\pi }{R_{2g+2}^{1/2}},$$
where the constants $`\underset{ยฏ}{c}(.)`$ are given by
$$c_j(k)=C_{jk}^1,C_{jk}=_{a_k}\frac{\pi ^{j1}d\pi }{R_{2g+2}^{1/2}}=2_{E_{2k1}}^{E_{2k}}\frac{z^{j1}dz}{R_{2g+2}^{1/2}(z)}.$$
The differentials fulfill
(2.6)
$$_{a_j}\zeta _k=\delta _{j,k},_{b_j}\zeta _k=\tau _{j,k},\tau _{j,k}=\tau _{k,j},1j,kg.$$
Now pick $`g`$ numbers (the Dirichlet eigenvalues)
(2.7)
$$(\widehat{\mu }_j)_{j=1}^g=(\mu _j,\sigma _j)_{j=1}^g$$
whose projections lie in the spectral gaps, that is, $`\mu _j[E_{2j1},E_{2j}]`$. Associated with these numbers is the divisor $`๐_{\underset{ยฏ}{\overset{^}{\mu }}}`$ which is one at the points $`\widehat{\mu }_j`$ and zero else. Using this divisor we introduce
(2.8)
$$\underset{ยฏ}{z}(p,n)=\underset{ยฏ}{\overset{^}{A}}_{p_0}(p)\underset{ยฏ}{\alpha }_{p_0}(๐_{\underset{ยฏ}{\overset{^}{\mu }}})n\underset{ยฏ}{A}_{\mathrm{}_{}}(\mathrm{}_+)\underset{ยฏ}{\overset{^}{\mathrm{\Xi }}}_{p_0}^g,\underset{ยฏ}{z}(n)=\underset{ยฏ}{z}(\mathrm{}_+,n),$$
where $`\underset{ยฏ}{\mathrm{\Xi }}_{p_0}`$ is the vector of Riemann constants
(2.9)
$$\widehat{\mathrm{\Xi }}_{p_0,j}=\frac{1_{k=1}^g\tau _{j,k}}{2},p_0=(E_0,0),$$
and $`\underset{ยฏ}{A}_{p_0}`$ ($`\underset{ยฏ}{\alpha }_{p_0}`$) is Abelโs map (for divisors). The hat indicates that we regard it as a (single-valued) map from $`\widehat{M}`$ (the fundamental polygon associated with $`๐`$) to $`^g`$. We recall that the function $`\theta (\underset{ยฏ}{z}(p,n))`$ has precisely $`g`$ zeros $`\widehat{\mu }_j(n)`$ (with $`\widehat{\mu }_j(0)=\widehat{\mu }_j`$), where $`\theta (\underset{ยฏ}{z})`$ is the Riemann theta function of $`๐`$.
Then our Jacobi operator $`H_q`$ is given by
$`a(n)^2`$ $`=`$ $`\stackrel{~}{a}^2{\displaystyle \frac{\theta (\underset{ยฏ}{z}(n+1))\theta (\underset{ยฏ}{z}(n1))}{\theta (\underset{ยฏ}{z}(n))^2}},`$
(2.10) $`b(n)`$ $`=`$ $`\stackrel{~}{b}+{\displaystyle \underset{j=1}{\overset{g}{}}}c_j(g){\displaystyle \frac{}{w_j}}\mathrm{ln}\left({\displaystyle \frac{\theta (\underset{ยฏ}{w}+\underset{ยฏ}{z}(n))}{\theta (\underset{ยฏ}{w}+\underset{ยฏ}{z}(n1))}}\right)|_{\underset{ยฏ}{w}=0}.`$
The constants $`\stackrel{~}{a}`$, $`\stackrel{~}{b}`$ depend only on the Riemann surface and will be defined in the next section.
It is well known that the spectrum of $`H_q`$ is purely absolutely continuous and consists of $`g+1`$ bands
(2.11)
$$\sigma (H_q)=\underset{j=0}{\overset{g}{}}[E_{2j},E_{2j+1}].$$
For further information and proofs we refer to , Section 9.
## 3. The Baker-Akhiezer function and the quasi-momentum map
The Baker-Akhiezer function $`\psi _q(p,n)=\psi _q(p,n,0)`$ is given by
(3.1)
$$\psi _q(p,n,n_0)=\sqrt{\frac{\theta (\underset{ยฏ}{z}(n_01))\theta (\underset{ยฏ}{z}(n_0))}{\theta (\underset{ยฏ}{z}(n1))\theta (\underset{ยฏ}{z}(n))}}\frac{\theta (\underset{ยฏ}{z}(p,n))}{\theta (\underset{ยฏ}{z}(p,n_0))}\mathrm{exp}\left((nn_0)_{p_0}^p\widehat{\omega }_{\mathrm{}_+,\mathrm{}_{}}\right),$$
where $`\omega _{\mathrm{}_+,\mathrm{}_{}}`$ is the normalized Abelian differential of the third kind with simple poles at $`\mathrm{}_\pm `$ and residues $`\pm 1`$, respectively. They are normalized such that $`\psi _q(p,n_0,n_0)=1`$.
The two branches
(3.2)
$$\psi _{q,\pm }(z,n)=\underset{j=0}{\overset{n1}{}}\varphi _{q,\pm }(z,j),$$
where (, (8.87))
(3.3)
$$\varphi _{q,\pm }(z,n)=\frac{1}{2a_q(n)}\left(zb_q(n)+\underset{j=1}{\overset{g}{}}\frac{\widehat{R}_j(n)}{z\mu _j(n)}\pm \frac{R_{2g+2}^{1/2}(z)}{_{j=1}^g(z\mu _j(n))}\right),$$
$$R_j(n)=\frac{R_{2g+1}^{1/2}(\mu _j(n))}{_{kj}(\mu _j(n)\mu _k(n))},\widehat{R}_j(n)=\sigma _j(n)R_j(n),$$
of the Baker-Akhiezer function are solutions of $`\tau _qu=zu`$, $`z`$, where $`\tau _q`$ is the difference expression associated with $`H_q`$. However, the Wronskian
(3.4)
$$W(\psi _{q,}(z),\psi _{q,+}(z))=\frac{R_{2g+2}^{1/2}(z)}{_{j=1}^g(z\mu _j)}$$
($`\mu _j=\mu _j(0)`$) shows that they are linearly dependent at the band edge $`E_j`$, $`0j2g+1`$.
The branch $`\psi _{q,\sigma _j}(z,n)`$ has a first order pole at $`\mu _j`$ if $`\mu _j`$ is away from the band edges
(3.5)
$$\underset{z\mu _j}{lim}(z\mu _j)\psi _{q,\sigma _j}(z,n)=\psi _{q,\sigma _j}(\mu _j,n,1)\frac{\widehat{R}_j(0)}{a_q(0)}$$
(use (3.3) and $`\psi _{q,\pm }(z,n)=\psi _{q,\pm }(z,n,1)\varphi _{q,\pm }(z,0)`$) and both branches have a square root singularity if $`\mu _j`$ coincides with a band edge $`E_l`$
(3.6)
$$\underset{z\mu _j}{lim}\sqrt{z\mu _j}\psi _{q,\pm }(z,n)=\pm \frac{\mathrm{i}^l_{kl}\sqrt{|E_lE_k|}}{2a_q(0)_{kj}\sqrt{E_l\mu _k}}\psi _{q,+}(E_l,n,1)$$
###### Lemma 3.1.
The solutions of $`\tau _qu=zu`$ can be characterized as follows.
(i) If $`R_{2g+2}(z)0`$, there exist two solutions satisfying
(3.7)
$$\psi _{q,\pm }(z,n)=\theta _\pm (z,n)w(z)^{\pm n},w(z)=\mathrm{exp}\left(_{p_0}^{(z,+)}\widehat{\omega }_{\mathrm{}_+,\mathrm{}_{}}\right),$$
with $`\theta _\pm (z,n)`$ quasi-periodic.
(ii) If $`R_{2g+2}(z)=0`$, $`z=E_l`$, there are two solutions satisfying
(3.8)
$$\psi _q(E_l,n)=\psi _{q,+}(E_l,n)=\psi _{q,}(E_l,n),\widehat{\psi }_q(E_l,n)=\psi _q(E_l,n)(\widehat{\theta }_l(n)+n),$$
where $`\widehat{\theta }_l(n)`$ is quasi-periodic.
###### Proof.
(ii). We construct a second linearly independent solution at $`z=E=E_l`$ using (see , (1.50))
(3.9)
$$s_q(E,n)=\underset{zE}{lim}a_q(0)\frac{\psi _{q,+}(z,n)\psi _{q,}(z,n)}{W(\psi _{q,}(z),\psi _{q,+}(z))},$$
where $`s_q(z,n)`$ denotes the fundamental solution of $`\tau _qu=zu`$ with initial conditions $`s_q(z,0)=0`$, $`s_q(z,1)=1`$. W.l.o.g. we assume that $`E_l`$ does not coincide with one of the Dirichlet eigenvalues $`\mu _j`$ (otherwise shift the base point). To derive an expression for $`\psi _{q,\pm }(z)`$ at $`z=E+ฯต^2`$ we start with
$$R_{2g+2}^{1/2}(z)=ฯต(\stackrel{~}{R}+O(ฯต^2)),\stackrel{~}{R}=\underset{jl}{}\sqrt{EE_j}.$$
Moreover,
$$W(\psi _{q,}(z),\psi _{q,+}(z))=\frac{\stackrel{~}{R}}{_{j=1}^g(E\mu _j)}ฯต(1+O(ฯต^2))$$
and for $`p=(E+ฯต^2,\pm )`$ (see (3.11) below),
$$_{p_0}^p\widehat{\omega }_{\mathrm{}_+,\mathrm{}_{}}=_{p_0}^E\widehat{\omega }\pm \beta ฯต+O(ฯต^3),\beta =\frac{2_{j=1}^g(E\lambda _j)}{\stackrel{~}{R}},$$
$$\underset{ยฏ}{z}(p,n)=\underset{ยฏ}{z}(E,n)\pm \underset{ยฏ}{\gamma }ฯต+O(ฯต^3),\underset{ยฏ}{\gamma }=\underset{j=1}{\overset{g}{}}\underset{ยฏ}{c}(j)\frac{2E^{j1}}{\stackrel{~}{R}},$$
and
$$\theta (\underset{ยฏ}{z}(p,n))=\theta (\underset{ยฏ}{z}(E,n))\pm \frac{\theta }{\underset{ยฏ}{z}}(\underset{ยฏ}{z}(E,n))\underset{ยฏ}{\gamma }ฯต+O(ฯต^3).$$
Using this to evaluate the limit $`\epsilon 0`$ shows
$$s_q(E,n)=2a_q(0)\underset{j=1}{\overset{g}{}}\frac{E\mu _j}{E\lambda _j}\widehat{\psi }_q(E,n)=\psi _q(E,n)(\widehat{\theta }(n)+n),$$
where
$$\widehat{\theta }(n)=\frac{1}{_{j=1}^g(E\lambda _j)}\underset{j,k=1}{\overset{g}{}}E^jc_k(j)\frac{}{w_k}\mathrm{ln}\theta (\underset{ยฏ}{z}(E,n)+\underset{ยฏ}{w}),$$
and finishes the proof. โ
###### Remark 3.2.
(i). Since $`\psi _q(z,n)`$ has a singularity if $`z=\mu _j`$ the solutions in Lemma 3.1 are not well-defined for those $`z`$. However, you can either remove the singularities of $`\psi _q(z,n)`$ or choose a different normalization point $`n_00`$ to see that solutions of the above type exist for every $`z`$.
(ii). In the periodic case Floquet theory tells you that there are two possible cases at a band edge: Either two (linearly independent) periodic solutions or one periodic and one linearly growing solution. The above lemma shows that the first case happens if the corresponding gap is closed and the second if the gap is open.
To understand the properties of $`\psi _{q,\pm }(z,n)`$ we need to investigate the quasi-momentum map
(3.10)
$$w(z)=\mathrm{exp}\left(_{p_0}^p\widehat{\omega }_{\mathrm{}_+,\mathrm{}_{}}\right),p=(z,+).$$
The differential $`\omega _{\mathrm{}_+,\mathrm{}_{}}`$ is given by
(3.11)
$$\omega _{\mathrm{}_+,\mathrm{}_{}}=\frac{_{j=1}^g(\pi \lambda _j)}{R_{2g+2}^{1/2}}d\pi ,$$
where the constants $`\lambda _j`$ have to be determined from the normalization
(3.12)
$$_{a_j}\omega _{\mathrm{}_+,\mathrm{}_{}}=2\underset{E_{2j1}}{\overset{E_{2j}}{}}\frac{_{j=1}^g(z\lambda _j)}{R_{2g+2}^{1/2}(z)}๐z=0,$$
which shows $`\lambda _j(E_{2j1},E_{2j})`$.
Since $`\lambda _j(E_{2j1},E_{2j})`$ the integrand is a Herglotz function and admits the following representation (c.f. , Appendix B)
(3.13)
$$\frac{_{j=1}^g(z\lambda _j)}{R_{2g+2}^{1/2}(z)}=_{\mathrm{}}^{\mathrm{}}\frac{1}{\lambda z}๐\stackrel{~}{\mu }(\lambda )$$
with the probability measure
(3.14)
$$d\stackrel{~}{\mu }(\lambda )=\frac{_{j=1}^g(\lambda \lambda _j)}{\pi \mathrm{i}R_{2g+2}^{1/2}(\lambda )}\chi _{\sigma (H_q)}(\lambda )d\lambda .$$
Hence
(3.15) $`g(z,\mathrm{})`$ $`=`$ $`{\displaystyle _{p_0}^p}\omega _{\mathrm{}_+,\mathrm{}_{}}={\displaystyle _{E_0}^z}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{1}{\lambda \zeta }}๐\stackrel{~}{\mu }(\lambda )๐\zeta `$
$`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}\mathrm{ln}\left({\displaystyle \frac{\lambda E_0}{\lambda z}}\right)๐\stackrel{~}{\mu }(\lambda ).`$
In particular, note that $`\text{Re}(g(z,\mathrm{}))`$ is the Greenโs function of the upper sheet $`\mathrm{\Pi }_+`$ with pole at $`\mathrm{}_+`$ and $`\stackrel{~}{\mu }`$ is the equilibrium measure of the spectrum (see , Thm. III.37). We will abbreviate $`g(z)=g(z,\mathrm{})`$.
The asymptotic expansion of $`\mathrm{exp}(g(z))`$ is given by (, (9.42))
(3.16)
$$\mathrm{exp}\left(_{p_0}^p\widehat{\omega }_{\mathrm{}_+,\mathrm{}_{}}\right)=\frac{\stackrel{~}{a}}{z}\left(1+\frac{\stackrel{~}{b}}{z}+O(\frac{1}{z^2})\right),z\mathrm{},$$
where $`\stackrel{~}{a}`$ is the capacity of the spectrum and
(3.17)
$$\stackrel{~}{b}=\frac{1}{2}\underset{j=0}{\overset{2g+2}{}}E_j\underset{j=1}{\overset{g}{}}\lambda _j.$$
###### Theorem 3.3.
The map $`g`$ is a bijection from the upper (resp. lower) half plane $`^\pm =\{z\pm \mathrm{Im}(z)>0\}`$ to
(3.18)
$$S^\pm =\{z\pm \mathrm{Re}(z)<0,0<\mathrm{Im}(z)<\pi \}\backslash \underset{j=1}{\overset{g}{}}[g(\lambda _j),g(E_{2j+1})]$$
such that $`\sigma (H_q)=\{z\mathrm{Re}(z)=0\}`$.
###### Proof.
By the Herglotz property of its integrand, the function $`g(z,\mathrm{})`$ satisfies the conditions of , Theorem 1(b) in Chapter VI, which shows that it is one-to-one.
To prove that $`g(z,\mathrm{})`$ is surjective, it suffices to show that the boundary of $`^+`$ is mapped to the boundary of $`S^+`$. Note that $`g(\lambda )`$ is negative for $`\lambda <E_0`$ and purely imaginary for $`\lambda [E_0,E_1]`$. At $`E_1`$, the real part starts to decrease from zero until it hits its minimum at $`\lambda _1`$ and increases again until it becomes 0 at $`E_2`$ (since all $`a`$-periods are zero), while the imaginary part remains constant. Proceeding like this we move along the boundary of $`S^+`$ as $`\lambda `$ moves along the real line. For $`\lambda >E_{2g+1}`$, $`g(\lambda )`$ is again negative. โ
###### Remark 3.4.
In the special case where $`H_q`$ is periodic the quasi-momentum is given by $`w(z)=\mathrm{exp}(\mathrm{i}N^1\mathrm{arccos}\mathrm{\Delta }(z))`$, where $`\mathrm{\Delta }(z)`$ is the Floquet discriminant, and our result is due to .
Therefore the map
$`w:^\pm `$ $``$ $`W^\pm =\{w|w|<1,\pm \text{Im}(w)>0\}\backslash {\displaystyle \underset{j=1}{\overset{g}{}}}[w(\lambda _j),w(E_{2j+1})]`$
(3.19) $`z`$ $``$ $`\mathrm{exp}(g(z))`$
is bijective. Denote $`W=W^+W^{}(1,1)`$, $`W_0=W\backslash \{0\}`$. If we identify corresponding points on the slits $`[w(\lambda _j),w(E_{2j+1})]`$ we obtain a Riemann surface $`๐`$ which is isomorphic to the upper sheet $`\mathrm{\Pi }_+`$.
###### Remark 3.5.
In the largest band edge $`E_{2g+1}`$ is chosen for $`p_0`$ and $`w`$ will map $`^\pm W^{}`$ in this case. Moreover, in the periodic case the slits $`[w(\lambda _j),w(E_{2j+1})]`$ appear at equal angles $`\frac{2\pi }{N}`$, where $`N`$ is the period.
Since $`zw(z)=\mathrm{exp}(g(z))`$ is a bijection, we consider the functions $`\psi _{q,\pm }`$ as functions of the new parameter $`w`$ whenever convenient. For notational simplicity we will write $`\psi _{q,\pm }(w,n)`$ for $`\psi _{q,\pm }(\lambda (w),n)`$ and similarly for other quantities. The functions $`\psi _{q,\pm }(w,n)`$ are meromorphic in $`๐`$ and continuous up to the boundary with the only possible singularities at the images of the Dirichlet eigenvalues $`w(\mu _j)`$ and at $`0`$. More precisely, denote by $`M_\pm `$ the sets of poles (and square root singularities if $`\mu _j=E_l`$) of the Weyl $`m`$-functions $`\stackrel{~}{m}_\pm (\lambda )`$, i.e. $`M_+M_{}=\{\mu _j\}_{j=1}^g`$ (see (3.2) and , Section 2.1). Note that $`\mu _jM_+M_{}`$ if and only if $`\mu _j=E_l`$. Then
* $`\psi _{q,\pm }(w,n)`$ are holomorphic in $`๐\backslash (\{w(\mu _j)\}_{j=1}^g\{0\})`$ and continuous on $`W\backslash \{w(\mu _j)\}`$.
* $`\psi _{q,\pm }(w,n)`$ has a simple pole at $`w(\mu _j)`$ if $`\mu _jM_\pm \backslash \{E_l\}`$, no pole if $`\mu _jM_\pm `$, and if $`\mu _j=E_l`$,
$$\psi _{q,\pm }(w,n)=\pm \frac{\mathrm{i}^lC(n)}{ww_l}+O(1),$$
where $`C(n)`$ is bounded and real.
* $`\overline{\psi _{q,\pm }(w,n)}=\psi _{q,}(w,n)`$ for $`|w|=1`$.
* At $`w=0`$ the following asymptotics hold
$$\psi _{q,\pm }(w,n)=(1)^n\left(\frac{_{m=0}^{n1}a_q(m)}{\stackrel{~}{a}^n}\right)^{\pm 1}w^{\pm n}(1+O(w)).$$
By Section 2.5 of the vector valued functions
(3.20)
$$\underset{ยฏ}{U}(\lambda ,n)=\sqrt{\frac{1}{4a_q(0)^2\pi \text{Im}(\stackrel{~}{m}_+(\lambda ))}}\left(\begin{array}{c}\psi _{q,+}(\lambda ,n)\\ \psi _{q,}(\lambda ,n)\end{array}\right)$$
form an orthonormal basis for the Hilbert space $`L^2(\sigma (H_q),^2,d\lambda )`$. The Weyl $`m`$-functions $`\stackrel{~}{m}_\pm (z)`$ satisfy (see , eq. (8.27))
(3.21)
$$\mathrm{Im}(\stackrel{~}{m}_\pm (\lambda ))=\frac{R_{2g+2}^{1/2}(\lambda )}{2\mathrm{i}a_q(0)^2_{j=1}^g(\lambda \mu _j)},\lambda \sigma (H_q).$$
Using our map $`w(z)=\mathrm{exp}(_{p_0}^{(z,+)}\widehat{\omega }_{\mathrm{}_+,\mathrm{}_{}})`$ we can transform this into an orthogonal basis on the unit circle.
###### Lemma 3.6.
Both functions $`\psi _{q,+}(w,n)`$ and $`\psi _{q,}(w,n)`$ form orthonormal bases in the Hilbert space $`L^2(S^1,\frac{1}{2\pi \mathrm{i}}d\omega )`$, where
(3.22)
$$d\omega (w)=\underset{j=1}{\overset{g}{}}\frac{\lambda (w)\mu _j}{\lambda (w)\lambda _j}\frac{dw}{w}.$$
###### Proof.
Just use
(3.23)
$$\frac{dw}{dz}=w\frac{_{j=1}^g(z\lambda _j)}{R_{2g+2}^{1/2}(z)}.$$
Observe that $`d\omega `$ is meromorphic on $`๐`$ with a simple pole at $`w=0`$. In particular, there are no poles at $`w(\lambda _j)`$.
###### Remark 3.7.
In the periodic case we have
(3.24)
$$\psi _{p,\pm }(\lambda )_N^2:=\underset{n=1}{\overset{N}{}}|\psi _{p,\pm }(\lambda ,n)|^2=N\underset{j=1}{\overset{N1}{}}\frac{\lambda \lambda _j}{\lambda \mu _j}.$$
## 4. Existence of Jost solutions
After we have these preparations out of our way, we come to the study of short-range perturbations $`H`$ of $`H_q`$ associated with sequences $`a`$, $`b`$ satisfying $`a(n)a_q(n)`$ and $`b(n)b_q(n)`$ as $`|n|\mathrm{}`$. More precisely, we will make the following assumption throughout this paper.
###### Hypothesis H. 4.1.
Let $`H`$ be a perturbation such that
(4.1)
$$\underset{n}{}|n|\left(|a(n)a_q(n)|+|b(n)b_q(n)|\right)<\mathrm{}.$$
We first establish existence of Jost solutions, that is solutions of the perturbed operator which asymptotically look like the Baker-Akhiezer solutions.
###### Theorem 4.2.
Assume (H.4.1). Then there exist solutions $`\psi _\pm (z,.)`$, $`z`$, of $`\tau \psi =z\psi `$ satisfying
(4.2)
$$\underset{n\pm \mathrm{}}{lim}|w(z)^n(\psi _\pm (z,n)\psi _{q,\pm }(z,n))|=0,$$
where $`\psi _{q,\pm }(z,.)`$ are the Baker-Akhiezer functions. Moreover, $`\psi _\pm (z,.)`$ are continuous (resp. holomorphic) with respect to $`z`$ whenever $`\psi _{q,\pm }(z,.)`$ are and inherit the properties $`(B1)`$ and $`(B2)`$, where now $`\psi _\pm (z,n)=\frac{\mathrm{i}^lC_\pm (n)}{\sqrt{z\mu _j}}+O(1)`$ and $`(B4)`$ has to be replaced by
(4.3)
$$\psi _\pm (z,n)=A_\pm (0)\left(\frac{_{m=0}^{n1}a(m)}{z^n}\right)^{\pm 1}(1+(B_\pm (0)\pm \underset{j=1}{\overset{n}{}}b(j{\scriptscriptstyle \genfrac{}{}{0pt}{}{0}{1}}))\frac{1}{z}+O(\frac{1}{z^2})),$$
where
$`A_+(n)`$ $`=`$ $`{\displaystyle \underset{j=n}{\overset{\mathrm{}}{}}}{\displaystyle \frac{a_q(j)}{a(j)}},B_+(n)={\displaystyle \underset{m=n+1}{\overset{\mathrm{}}{}}}(b_q(m)b(m)),`$
(4.4) $`A_{}(n)`$ $`=`$ $`{\displaystyle \underset{j=\mathrm{}}{\overset{n1}{}}}{\displaystyle \frac{a_q(j)}{a(j)}},B_{}(n)={\displaystyle \underset{m=\mathrm{}}{\overset{n1}{}}}(b_q(m)b(m)).`$
###### Proof.
The proof can be done as in the periodic case (see e.g., , or , Section 7.5). The only problem is to show that the second solution at a band edge grows at most linearly. In the periodic case this follows from Floquet theory, here we just use Lemma 3.1. โ
From this result we obtain a complete characterization of the spectrum of $`H`$.
###### Theorem 4.3.
Assume (H.4.1). Then we have $`\sigma _{ess}(H)=\sigma (H_q)`$, the point spectrum of $`H`$ is finite and confined to the spectral gaps of $`H_q`$, that is, $`\sigma _p(H)\backslash \sigma (H_q)`$. Furthermore, the essential spectrum of $`H`$ is purely absolutely continuous.
###### Proof.
Again the proof can be done as in the periodic case (see e.g., or , Section 7.5). โ
## 5. The transformation operator
We define the kernel of the transformation operator as the Fourier coefficients of the Jost solutions $`\psi _\pm (w,n)`$ with respect to the orthonormal system given in Lemma 3.6, $`\{\psi _{q,\pm }(w,n)\}_n`$,
(5.1)
$$K_\pm (n,m):=\frac{1}{2\pi \mathrm{i}}_{|w|=1}\psi _\pm (w,n)\psi _{q,}(w,m)๐\omega (w).$$
By the Cauchy theorem, this integral equals the residue at $`w=0`$,
(5.2)
$$K_\pm (n,m)=\mathrm{Res}_0\frac{1}{w}\psi _\pm (w,n)\psi _{q,}(w,m).$$
In particular, since $`\psi _\pm (w,n)\psi _{q,}(w,m)=O(w^{\pm (nm)})`$, we conclude
(5.3)
$$K_\pm (n,m)=0,\pm (mn)<0.$$
###### Lemma 5.1.
Assume H.4.1. The Jost solutions $`\psi _\pm (w,n)`$ can be represented as
(5.4)
$$\psi _\pm (w,n)=\underset{m=n}{\overset{\pm \mathrm{}}{}}K_\pm (n,m)\psi _{q,\pm }(w,m),|w|=1,$$
where the kernels $`K_\pm (n,.)`$ satisfy $`K_\pm (n,m)=0`$ for $`\pm m<\pm n`$ and
(5.5)
$$|K_\pm (n,m)|C\underset{j=[\frac{m+n}{2}]\pm 1}{\overset{\pm \mathrm{}}{}}\left(|a(j)a_q(j)|+|b(j)b_q(j)|\right),\pm m>\pm n.$$
The constant $`C`$ depends only on $`H_q`$ and the value of the sum in (4.1).
###### Proof.
We prove the estimate for $`K_+(n,m)`$ and omit โ$`+`$โ and โ$`z`$โ whenever possible. Define $`\phi (n)=\psi (n)K(n,n)^1`$, then $`\phi `$ fulfills
(5.6)
$$\phi (n)=\psi _q(n)+\underset{m=n+1}{\overset{\mathrm{}}{}}J(n,m)\phi (m),$$
where
(5.7)
$$J(z,n,m)=\stackrel{~}{a}(m1)\frac{s_q(z,n,m1)}{a_q(m1)}+\stackrel{~}{b}(m)\frac{s_q(z,n,m)}{a_q(m)}$$
with the abbreviation
(5.8)
$$\stackrel{~}{a}(m)=\frac{a(m)^2}{a_q(m)}a_q(m),\stackrel{~}{b}(m)=b(m)b_q(m).$$
On the other hand, $`\phi (n)`$ is given by
$$\phi (n)=\underset{m=n}{\overset{\mathrm{}}{}}\kappa (n,m)\psi _q(m),\kappa (n,m)=\frac{K(n,m)}{K(n,n)},$$
therefore
(5.9)
$$\underset{m=n}{\overset{\mathrm{}}{}}\kappa (n,m)\psi _q(m)=\underset{m=n+1}{\overset{\mathrm{}}{}}J(n,m)\psi _q(m)+\underset{m=n+1}{\overset{\mathrm{}}{}}\underset{l=m+1}{\overset{\mathrm{}}{}}J(n,m)\kappa (m,l)\psi _q(l).$$
Multiplying both sides of (5.9) by $`\psi _{q,}(k)`$ and integrating over the unit circle yields
(5.10)
$$\kappa (n,k)=\underset{m=n+1}{\overset{\mathrm{}}{}}\mathrm{\Gamma }(n,m,m,k)+\underset{m=n+1}{\overset{\mathrm{}}{}}\underset{l=n+1}{\overset{\mathrm{}}{}}\mathrm{\Gamma }(n,m,l,k)\kappa (m,l),$$
where
(5.11)
$$\mathrm{\Gamma }(n,m,l,k)=\frac{1}{2\pi \mathrm{i}}_{|w|=1}J(w,n,m)\psi _{q,+}(w,l)\psi _{q,}(w,k)๐\omega (w).$$
Using , (1.50),
(5.12)
$$\frac{s_q(n,m)}{a(m)}=\frac{\psi _{q,+}(m)\psi _{q,}(n)\psi _{q,+}(n)\psi _{q,}(m)}{W(\psi _{q,+},\psi _{q,})},$$
we obtain
(5.13)
$$\mathrm{\Gamma }(n,m,l,k)=\stackrel{~}{b}(m)\mathrm{\Gamma }_q(n,m,l,k)+\stackrel{~}{a}(m)\mathrm{\Gamma }_q(n,m1,l,k)$$
with
$`\mathrm{\Gamma }_q(n,m,l,k)`$ $`=`$ $`\mathrm{\Gamma }_0(m,n,l,k)\mathrm{\Gamma }_0(n,m,l,k),`$
(5.14) $`\mathrm{\Gamma }_0(n,m,l,k)`$ $`=`$ $`{\displaystyle \frac{1}{2\pi \mathrm{i}}}{\displaystyle _{w(\gamma )}}{\displaystyle \frac{\psi _{q,+}(w,n)\psi _{q,}(w,m)\psi _{q,+}(w,l)\psi _{q,}(w,k)}{W(\psi _{q,+}(w),\psi _{q,}(w))}}๐\omega (w)`$
$`=`$ $`{\displaystyle \frac{1}{2\pi \mathrm{i}}}{\displaystyle _\gamma }{\displaystyle \frac{\psi _{q,+}(z,n)\psi _{q,}(z,m)\psi _{q,+}(z,l)\psi _{q,}(z,k)}{W(\psi _{q,+}(z),\psi _{q,}(z))}}{\displaystyle \frac{(z\mu _j)}{R_{2g+2}^{1/2}(z)}}๐z`$
$`=`$ $`{\displaystyle \frac{1}{2\pi \mathrm{i}}}{\displaystyle _\gamma }{\displaystyle \frac{\psi _{q,+}(z,n)\psi _{q,}(z,m)\psi _{q,+}(z,l)\psi _{q,}(z,k)}{W(\psi _{q,+}(z),\psi _{q,}(z))^2}}๐z.`$
Here $`\gamma `$ is a path on the upper sheet encircling the spectrum. The integrand of $`\mathrm{\Gamma }_0`$ is meromorphic on the Riemann surface $`๐`$ with poles of order one at $`E_j`$ and poles of order $`O(z^{\pm (nm+lk)2})`$ near $`\mathrm{}_\pm `$ (there are no poles at the Dirichlet eigenvalues $`\mu _j`$ ). We apply the residue theorem twice, first on the side of $`\gamma `$ including $`\mathrm{}_+`$, then on the other side including the spectrum (and thus $`\mathrm{}_{}`$)
(5.15) $`\mathrm{\Gamma }_0(n,m,l,k)`$ $`=`$ $`\mathrm{Res}_\mathrm{}_+{\displaystyle \frac{\psi _{q,+}(n)\psi _{q,}(m)\psi _{q,+}(l)\psi _{q,}(k)}{W(\psi _{q,+},\psi _{q,})^2}}`$
$`=`$ $`\left(\mathrm{Res}_{\mathrm{}_{}}+{\displaystyle \underset{j=0}{\overset{2g+1}{}}}\mathrm{Res}_{E_j}\right)\left({\displaystyle \frac{\psi _{q,+}(n)\psi _{q,}(m)\psi _{q,+}(l)\psi _{q,}(k)}{W(\psi _{q,+},\psi _{q,})^2}}\right).`$
The order of the poles at $`\mathrm{}_\pm `$ implies
$$\mathrm{\Gamma }_0(n,m,l,k)=\{\begin{array}{cc}\underset{j=0}{\overset{2g+1}{}}\mathrm{Res}_{E_j}\frac{\psi _{q,+}(n)\psi _{q,}(m)\psi _{q,+}(l)\psi _{q,}(k)}{W(\psi _{q,+},\psi _{q,})^2}& nm+lk<0\\ 0& nm+lk0,\end{array}$$
which shows that $`\mathrm{\Gamma }_0(n,m,l,k)`$ is real and bounded since $`\psi _{q,+}(E,.)=\psi _{q,}(E,.)`$ are (if $`\mu _j=E_l`$, use (B2)). Together with (5) this yields
(5.16)
$$\mathrm{\Gamma }_0(n,m,l,k)=\overline{\mathrm{\Gamma }_0(m,n,k,l)}=\mathrm{\Gamma }_0(m,n,k,l)=\mathrm{\Gamma }_0(n,m,k,l).$$
Moreover,
$`\mathrm{\Gamma }_q(n,m,l,k)`$ $`=`$ $`0,lk|mn|,`$
(5.17) $`\mathrm{\Gamma }_q(n,m,l,k)`$ $`=`$ $`\mathrm{\Gamma }_q(m,n,k,l)=\mathrm{\Gamma }_q(n,m,k,l),`$
which then implies
(5.18)
$$\mathrm{\Gamma }_q(n,m,l,k)=\{\begin{array}{cc}\text{sign}(nm)\underset{j=0}{\overset{2g+1}{}}\mathrm{Res}_{E_j}\frac{\psi _{q,+}(n)\psi _{q,}(m)\psi _{q,+}(l)\psi _{q,}(k)}{W(\psi _{q,+},\psi _{q,})^2}& |lk|<|mn|\\ 0& |lk||mn|\end{array}$$
and $`\mathrm{\Gamma }(n,m,l,k)=0`$ for $`|lk|mn`$ if $`m>n`$. Note that the residue at $`E_j`$ is given by
(5.19)
$$\frac{2_{\mathrm{}=1}^g(E_j\mu _{\mathrm{}})^2}{_\mathrm{}j(E_jE_{\mathrm{}})}\psi _q(E_j,n)\psi _q(E_j,m)\psi _q(E_j,l)\psi _q(E_j,k).$$
Now we obtain for $`\kappa (n,k)`$
(5.20) $`\kappa (n,k)`$ $`=`$ $`{\displaystyle \underset{m=n+1}{\overset{\mathrm{}}{}}}\mathrm{\Gamma }(n,m,m,k)+{\displaystyle \underset{m=n+1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{l=m+1}{\overset{\mathrm{}}{}}}\mathrm{\Gamma }(n,m,l,k)\kappa (m,l)`$
$`=`$ $`{\displaystyle \underset{m=[\frac{n+k}{2}]+1}{\overset{\mathrm{}}{}}}\mathrm{\Gamma }(n,m,m,k)+{\displaystyle \underset{m=n+1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{l=n+km+1}{\overset{m+kn1}{}}}\mathrm{\Gamma }(n,m,l,k)\kappa (m,l),`$
since $`\mathrm{\Gamma }(n,m,m,k)0`$ only if $`|mk|<mn`$ implying $`m>\frac{n+k}{2}`$. In the third sum of (5.20) we need that $`|m+\delta k|<mn`$ for $`\delta 1`$ which yields $`\delta <kn`$ and $`\delta >n+k2m`$. Two remarks might be in order: $`m+kn1n+km+1`$ since $`mnnm+2`$, and the starting point $`l=n+km+1`$ of the third sum actually has a lower limit, namely $`m\frac{n+k}{2}`$, since we require $`lm+1`$ for $`\kappa (m,l)0,1`$. Note that
$`\left|{\displaystyle \underset{m=[\frac{n+k}{2}]+1}{\overset{\mathrm{}}{}}}\mathrm{\Gamma }(n,m,m,k)\right|`$ $``$ $`D{\displaystyle \underset{m=[\frac{n+k}{2}]+1}{\overset{\mathrm{}}{}}}|\stackrel{~}{b}(m)+\stackrel{~}{a}(m)|=:\widehat{q}(\frac{n+k}{2}),`$
$`\left|{\displaystyle \underset{l=n+km+1}{\overset{m+kn1}{}}}|\mathrm{\Gamma }(n,m,l,k)|\right|`$ $``$ $`D(mn1)|\stackrel{~}{b}(m)+\stackrel{~}{a}(m)|=:\widehat{c}(m)\mathrm{}^1(),`$
where $`D`$ is the estimate provided by (5.18), (5.19). We set up the following iteration procedure
$`\kappa _0(n,k)`$ $`=`$ $`{\displaystyle \underset{m=[\frac{n+k}{2}]+1}{\overset{\mathrm{}}{}}}\mathrm{\Gamma }(n,m,m,k),`$
(5.21) $`\kappa _j(n,k)`$ $`=`$ $`{\displaystyle \underset{m=n+1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{l=n+km+1}{\overset{m+kn1}{}}}\mathrm{\Gamma }(n,m,l,k)\kappa _{j1}(m,l).`$
Then using induction one has
(5.22)
$$|\kappa _j(n,k)|\widehat{q}(\frac{n+k}{2})\frac{\left(_{m=n+1}^{\mathrm{}}\widehat{c}(m)\right)^j}{j!}$$
and hence the iteration converges and implies the estimate
(5.23)
$$|\kappa (n,k)|=\left|\underset{j=0}{\overset{\mathrm{}}{}}\kappa _j(n,k)\right|\widehat{q}(\frac{n+k}{2})\mathrm{exp}\left(\underset{m=n+1}{\overset{\mathrm{}}{}}\widehat{c}(m)\right).$$
Associated with $`K_\pm (n,m)`$ is the operator
(5.24)
$$(๐ฆ_\pm f)(n)=\underset{m=n}{\overset{\pm \mathrm{}}{}}K_\pm (n,m)f(m),f\mathrm{}_\pm ^{\mathrm{}}(,),$$
which acts as a transformation operator for the pair $`\tau `$, $`\tau _q`$.
###### Theorem 5.2.
Let $`\tau _q`$ and $`\tau `$ be the quasi-periodic and perturbed Jacobi difference expression, respectively. Then
(5.25)
$$\tau ๐ฆ_\pm f=๐ฆ_\pm \tau _qf,f\mathrm{}_\pm ^{\mathrm{}}(,).$$
###### Proof.
It suffices to show that $`HK_\pm =K_\pm H_q`$.
(5.26) $`HK_\pm (n,m)`$ $`=`$ $`{\displaystyle \frac{1}{2\pi \mathrm{i}}}{\displaystyle _{|w|=1}}H\psi _\pm (w,n)\psi _{q,}(w,m)๐\omega (w)`$
$`=`$ $`{\displaystyle \frac{1}{2\pi \mathrm{i}}}{\displaystyle _{|w|=1}}\lambda (w)\psi _\pm (w,n)\psi _{q,}(w,m)๐\omega (w)`$
$`=`$ $`{\displaystyle \frac{1}{2\pi \mathrm{i}}}{\displaystyle _{|w|=1}}\psi _\pm (w,n)H_q\psi _{q,}(w,m)๐\omega (w).`$
###### Lemma 5.3.
For $`n`$ we have
(5.27) $`{\displaystyle \frac{a(n)}{a_q(n)}}`$ $`=`$ $`{\displaystyle \frac{K_+(n+1,n+1)}{K_+(n,n)}}={\displaystyle \frac{K_{}(n,n)}{K_{}(n+1,n+1)}},`$
$`b(n)b_q(n)`$ $`=`$ $`a_q(n){\displaystyle \frac{K_+(n,n+1)}{K_+(n,n)}}a_q(n1){\displaystyle \frac{K_+(n1,n)}{K_+(n1,n1)}}`$
$`=`$ $`a_q(n1){\displaystyle \frac{K_{}(n,n1)}{K_{}(n,n)}}a_q(n){\displaystyle \frac{K_{}(n+1,n)}{K_{}(n+1,n+1)}},`$
###### Proof.
Consider the equation of the transformation operator $`H๐ฆ_\pm =๐ฆ_\pm H_q`$, which is equivalent to (c.f. (5.26))
$$a(n1)K_\pm (n1,m)+b(n)K_\pm (n,m)+a(n)K_\pm (n+1,m)=$$
$$=a_q(m1)K_\pm (n,m1)+b_q(m)K_\pm (n,m)+a_q(m)K_\pm (n,m+1).$$
Evaluating at $`m=n`$ we obtain the first equation and at $`m=n1`$ the second. โ
In particular, observe
(5.28)
$$K_\pm (n,n)=A_\pm (n),K_\pm (n,n\pm 1)=\frac{A_\pm (n)}{a_q(n\genfrac{}{}{0pt}{}{0}{1})}B_\pm (n).$$
## 6. The scattering matrix
Let $`H_q`$ be a given quasi-periodic Jacobi operator and $`H`$ a perturbation of $`H_q`$ satisfying Hypothesis H.4.1. To set up scattering theory for the pair $`(H,H_q)`$ we proceed as usual.
The Wronskian of our Jost functions can be evaluated as $`n\pm \mathrm{}`$ and is given by
(6.1)
$$W(\psi _\pm (\lambda ),\overline{\psi _\pm (\lambda )})=W_q(\psi _{q,\pm }(\lambda ),\psi _{q,}(\lambda ))=\frac{R_{2g+2}^{1/2}(\lambda )}{_{j=1}^g(\lambda \mu _j)},\lambda \sigma (H_q).$$
Hence $`\psi _\pm (\lambda )`$, $`\overline{\psi _\pm (\lambda )}`$ are linearly independent for $`\lambda `$ in the interior of $`\sigma (H_q)`$ and we consider the scattering relations
(6.2)
$$\psi _\pm (\lambda ,n)=\alpha (\lambda )\overline{\psi _{}(\lambda ,n)}+\beta _{}(\lambda )\psi _{}(\lambda ,n),\lambda \sigma (H_q),$$
where
(6.3) $`\alpha (\lambda )`$ $`=`$ $`{\displaystyle \frac{W(\psi _{}(\lambda ),\psi _\pm (\lambda ))}{W(\psi _{}(\lambda ),\overline{\psi _{}(\lambda )})}}={\displaystyle \frac{_{j=1}^g(\lambda \mu _j)}{R_{2g+2}^{1/2}(\lambda )}}W(\psi _{}(\lambda ),\psi _+(\lambda )),`$
$`\beta _\pm (\lambda )`$ $`=`$ $`{\displaystyle \frac{W(\psi _{}(\lambda ),\overline{\psi _\pm (\lambda )})}{W(\psi _\pm (\lambda ),\overline{\psi _\pm (\lambda )})}}={\displaystyle \frac{_{j=1}^g(\lambda \mu _j)}{R_{2g+2}^{1/2}(\lambda )}}W(\psi _{}(\lambda ),\overline{\psi _\pm (\lambda )}).`$
While $`\alpha (\lambda )`$ is only defined for $`\lambda \sigma (H_q)`$, (6.3) may be used as a definition for $`\lambda \backslash \{E_j\}`$. Therefore $`\alpha (w)`$ can be continued as a holomorphic function on $`๐`$ and it is continuous up to the boundary except possibly at the band edges.
###### Remark 6.1.
Note that $`\alpha (\lambda )`$ does not depend on the normalization of $`\psi _\pm (\lambda )`$ at the base point $`n_0=0`$ whereas $`\beta _\pm =\beta _{\pm ,0}`$ does. Using $`\psi _\pm (z,n,n_0)=\psi _{q,\pm }(z,n_0)^1\psi _\pm (z,n)`$ and
$$W((\psi _+(\lambda ),\psi _{}(\lambda ))=\underset{j=1}{\overset{g}{}}\frac{\lambda \mu _j(n_0)}{\lambda \mu _j}W((\psi _+(\lambda ,.,n_0),\psi _{}(\lambda ,.,n_0))$$
we see
(6.4)
$$\beta _{\pm ,0}(\lambda )=\frac{\psi _{q,}(\lambda ,n_0)}{\psi _{q,\pm }(\lambda ,n_0)}\beta _{\pm ,n_0}(\lambda ).$$
A direct calculation shows
(6.5)
$$\alpha (\overline{w})=\overline{\alpha (w)},\beta _\pm (\overline{w})=\overline{\beta _\pm (w)}=\beta _{}(w)$$
and the Plรผcker identity (c.f. , (2.169)) implies
(6.6)
$$|\alpha (w)|^2=1+|\beta _\pm (w)|^2,|w|=1.$$
We will denote the eigenvalues of $`H`$ by
(6.7)
$$\sigma _p(H)=\{\rho _j\}_{j=1}^q.$$
Our next aim is to study the behavior of $`\alpha (\lambda )`$ at the eigenvalues $`\rho _j`$, therefore we modify the Jost solutions $`\psi _\pm (\lambda ,n)`$ according to their poles at $`\mu _j`$ and define the following eigenfunctions $`\widehat{\psi }_\pm (\lambda ,.)`$
(6.8)
$$\widehat{\psi }_+(\lambda ,.)=\underset{\mu _lM_+}{}(\lambda \mu _l)\psi _+(\lambda ,.),\widehat{\psi }_{}(\lambda ,.)=\underset{\mu _lM_{}\backslash \{E_j\}}{}(\lambda \mu _l)\psi _{}(\lambda ,.).$$
Define $`\widehat{\psi }_{q,\pm }(\lambda ,.)`$ accordingly. Moreover, $`\widehat{\psi }_\pm (\rho _j,n)=c_j^\pm \widehat{\psi }_{}(\rho _j,n)`$ with $`c_j^+c_j^{}=1`$. The norming constants $`\gamma _{\pm ,j}`$ are defined by
(6.9)
$$\frac{1}{\gamma _{\pm ,j}}:=\underset{m}{}|\widehat{\psi }_\pm (\rho _j,m)|^2.$$
To compute the derivative of $`\alpha (\lambda )`$ at $`\rho _j`$, note that
(6.10)
$$\alpha (\lambda )=\frac{W(\widehat{\psi }_{}(\lambda ),\widehat{\psi }_+(\lambda ))}{R_{2g+2}^{1/2}(\lambda )}.$$
By virtue of , Lemma 2.4,
(6.11)
$$\frac{d}{d\lambda }W(\widehat{\psi }_{}(\lambda ),\widehat{\psi }_+(\lambda ))|_{\rho _j}=\underset{k}{}\widehat{\psi }_{}(\rho _j,k)\widehat{\psi }_+(\rho _j,k)=\frac{1}{c_j^\pm \gamma _{\pm ,j}}.$$
Therefore
(6.12)
$$\frac{d}{d\lambda }\alpha (\lambda )|_{\rho _j}=\frac{W^{}(\widehat{\psi }_{}(\rho _j),\widehat{\psi }_+(\rho _j))}{R_{2g+2}^{1/2}(\rho _j)}=\frac{1}{c_j^\pm \gamma _{\pm ,j}R_{2g+2}^{1/2}(\rho _j)}.$$
From (6.12) we obtain a connection between the left and right norming constants
(6.13)
$$\gamma _{+,j}\gamma _{,j}=\frac{1}{(\alpha ^{}(\rho _j))^2R_{2g+2}(\rho _j)}.$$
As a last preparation, we study the behavior of $`\alpha (w)`$ as $`w0`$. By (4.3),
(6.14)
$$W(\psi _{}(w),\psi _+(w))=A_{}(0)A_+(0)\stackrel{~}{a}w^1+O(w)$$
and
(6.15)
$$\frac{R_{2g+2}^{1/2}(\lambda (w))}{_{j=1}^g(\lambda (w)\lambda _j)}=\stackrel{~}{a}w^1+O(1),$$
therefore $`\alpha ^1(w)`$ is bounded at $`0`$ with
(6.16)
$$\alpha (0)=\underset{j=\mathrm{}}{\overset{\mathrm{}}{}}\frac{a_q(j)}{a(j)}.$$
We now define the scattering matrix
(6.17)
$$S(w)=\left(\begin{array}{cc}T(w)& R_{}(w)\\ R_+(w)& T(w)\end{array}\right),|w|=1,$$
where $`T(w):=\alpha ^1(w)`$ and $`R_\pm (w):=\alpha ^1(w)\beta _\pm (w)`$ are called transmission and reflection coefficients. Equations (6.5) and (6.6) imply
###### Lemma 6.2.
The scattering matrix $`S(w)`$ is unitary. The coefficients $`T(w)`$, $`R_\pm (w)`$ are bounded for $`|w|=1`$, continuous for $`|w|=1`$ except at possibly $`w_l=w(E_l)`$, fulfill
(6.18) $`|T(w)|^2+|R_\pm (w)|^2`$ $`=`$ $`1,|w|=1,`$
(6.19) $`T(w)R_+(\overline{w})+T(\overline{w})R_{}(w)`$ $`=`$ $`0,|w|=1`$
and $`\overline{T(w)}=T(\overline{w})`$, $`\overline{R_\pm (w)}=R_\pm (\overline{w})`$ for $`|w|=1`$
Moreover, $`R_{2g+2}^{1/2}(w)T(w)^1`$ is continuous (in particular $`T(w)`$ can only vanish at $`w_l`$) and
(6.20)
$$\begin{array}{cc}\underset{ww_l}{lim}R_{2g+2}^{1/2}(w)\frac{R_\pm (w)+1}{T(w)}=0,\hfill & w_lw(\mu _j)\hfill \\ \underset{ww_l}{lim}R_{2g+2}^{1/2}(w)\frac{R_\pm (w)1}{T(w)}=0,\hfill & w_l=w(\mu _j)\hfill \end{array}.$$
The transmission coefficient $`T(w)`$ has a meromorphic continuation to $`๐`$ with simple poles at $`w(\rho _j)`$,
(6.21)
$$\left(\mathrm{Res}_{\rho _j}T(\lambda )\right)^2=\gamma _{+,j}\gamma _{,j}R_{2g+2}(\rho _j).$$
In addition, $`T(z)`$ as $`z\backslash \sigma (H_q)`$ and
(6.22)
$$T(0)=\frac{1}{K_+(n,n)K_{}(n,n)}=\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}\frac{a(m)}{a_q(m)},$$
where $`K_\pm (n,n)`$ are the coefficients of the transformation operators.
###### Proof.
To show (6.20) we use the definition (6.3),
$$R_{2g+2}^{1/2}(\lambda )\frac{R_\pm (\lambda )+1}{T(\lambda )}=\underset{j=1}{\overset{g}{}}(\lambda \mu _j)\left(W(\psi _{}(\lambda ),\psi _+(\lambda ))W(\psi _{}(\lambda ),\overline{\psi _\pm (\lambda )})\right).$$
There are two cases to distinguish: If $`\mu _jE_l`$ then $`\psi _\pm `$ are continuous and real at $`\lambda =E_l`$ and the two Wronskians cancel. Otherwise, if $`\mu _j=E_l`$ they are purely imaginary (by property (B2) of the Jost functions) and the two terms are equal in the limit and add up. โ
The sets
(6.23)
$$S_\pm (H)=\{R_\pm (w),|w|=1;(\rho _j,\gamma _{\pm ,j}),1jq\}$$
are called left/right scattering data for $`H`$.
First we want to show that the transmission coefficient can be reconstructed from either left or right scattering data.
Let $`g(w,w_0)`$ be the Green function associated with $`๐`$ and let
(6.24)
$$\mu (w,w_0)dw_0=\frac{g}{r}(w,r\mathrm{e}^{\mathrm{i}\theta })|_{r=1^{}}\mathrm{e}^{\mathrm{i}\theta }d\theta ,w_0=\mathrm{e}^{\mathrm{i}\theta },$$
be the corresponding harmonic measure on the boundary (see, e.g., ). Since $`W_0`$ is simply connected, we can choose a function $`h(w,v)`$ such that $`\widehat{g}(w,w_0)=g(w,w_0)+\mathrm{i}h(w,w_0)`$ is analytic in $`W_0`$. Clearly $`\widehat{g}`$ is only well-defined up to an imaginary constant and it will not be analytic on $`๐\backslash \{0\}`$ in general. Similarly we can find a corresponding $`\nu (w,w_0)`$ and set $`\widehat{\mu }(w,w_0)=\mu (w,w_0)+\mathrm{i}\nu (w,w_0)`$.
###### Theorem 6.3.
Either one of the sets $`S_\pm (H)`$ determines the other and $`T(w)`$ via the Poisson-Jensen type formula
(6.25)
$$T(w)=\mathrm{exp}\left(\underset{j=1}{\overset{q}{}}\widehat{g}(w,w(\rho _j))\right)\mathrm{exp}\left(\frac{1}{2}_{|w|=1}\mathrm{ln}(1|R_\pm (w_0)|^2)\widehat{\mu }(w,w_0)๐w_0\right),$$
where the constant of $`\widehat{g}`$ has to be chosen such that $`T(0)>0`$, and
$$\frac{R_{}(w)}{R_+(\overline{w})}=\frac{T(w)}{T(\overline{w})},\gamma _{+,j}\gamma _{,j}=\frac{\left(\mathrm{Res}_{\rho _j}T(\lambda )\right)^2}{_{l=0}^{2g+1}(\rho _jE_l)}.$$
###### Proof.
It suffices to prove the formula for $`T(w)`$, since evaluating the residua provides $`\gamma _{\pm ,j}`$, together with $`\{\lambda _l\}`$, $`\{E_l\}`$. The formula for $`T(w)`$ holds by , Theorem 1 at least when taking absolute values. Since both sides are analytic, and have equal absolute values, they can only differ by a constant of absolute value one. But both sides are positive at $`w=0`$ and hence this constant is one. โ
Note that neither the Blaschke factors nor the outer function in (6.25) are single valued on $`๐`$ in general. In particular, the eigenvalues cannot be chosen arbitrarily, which was first observed in .
## 7. The Gelโfand-Levitan-Marchenko equations
In this section we want to derive a procedure which allows the reconstruction of the Jacobi operator $`H`$ with asymptotically quasi-periodic coefficients from its scattering data $`S_\pm (H)`$. This will be achieved by deriving an equation for $`K_\pm (n,m)`$ which is generally known as Gelโfand-Levitan-Marchenko equation.
Since $`K_\pm (n,m)`$ are essentially the Fourier coefficients of the Jost solutions $`\psi _\pm (w,n)`$ we compute the Fourier coefficients of the scattering relations (6.2). Therefore we multiply
(7.1)
$$T(w)\psi _{}(w,n)=R_\pm (w)\psi _\pm (w,n)+\overline{\psi _\pm (w,n)}$$
by $`(2\pi \mathrm{i})^1\psi _{q,\pm }(w,m)d\omega `$, where $`\pm m\pm n`$, and integrate around the unit circle. First we evaluate the right hand side of (7.1) using (5.1)
(7.2) $`{\displaystyle \frac{1}{2\pi \mathrm{i}}}{\displaystyle _{|w|=1}}\overline{\psi _+(w,n)}\psi _{q,+}(w,m)๐\omega (w)`$ $`=`$ $`K_+(n,m),`$
$`{\displaystyle \frac{1}{2\pi \mathrm{i}}}{\displaystyle _{|w|=1}}R_+(w)\psi _+(w,n)\psi _{q,+}(w,m)๐\omega (w)`$ $`=`$ $`{\displaystyle \underset{l=n}{\overset{\mathrm{}}{}}}K_+(n,l)\stackrel{~}{F}^+(l,m),`$
where
(7.3)
$$\stackrel{~}{F}^+(l,m)=\frac{1}{2\pi \mathrm{i}}_{|w|=1}R_+(w)\psi _{q,+}(w,l)\psi _{q,+}(w,m)๐\omega (w).$$
Note that $`\stackrel{~}{F}^+(l,m)=\stackrel{~}{F}^+(m,l)`$ is real.
To evaluate the left hand side of (7.1) we use the residue theorem. The only poles are at the eigenvalues and at $`0`$ if $`n=m`$, hence
$`{\displaystyle \frac{1}{2\pi \mathrm{i}}}{\displaystyle _{|w|=1}}T(w)\psi _{}(w,n)\psi _{q,+}(w,m)๐\omega (w)`$
$`={\displaystyle \frac{\delta (n,m)}{K_+(n,n)}}+{\displaystyle \underset{j=1}{\overset{q}{}}}\mathrm{Res}_{\rho _j}\left({\displaystyle \frac{T(\lambda )\widehat{\psi }_{}(\lambda ,n)\widehat{\psi }_{q,+}(\lambda ,m)}{R_{2g+2}^{1/2}(\lambda )}}\right).`$
Here $`\delta (n,m)`$ is one for $`m=n`$ and zero else. By (6.12) the residua at the eigenvalues are given by
(7.4)
$$\mathrm{Res}_{\rho _j}\left(\frac{T(\lambda )\widehat{\psi }_{}(\lambda ,n)\widehat{\psi }_{q,+}(\lambda ,m)}{R_{2g+2}^{1/2}(\lambda )}\right)=\gamma _{+,j}\widehat{\psi }_+(\rho _j,n)\widehat{\psi }_{q,+}(\rho _j,m).$$
Collecting all terms yields
(7.5)
$$K_\pm (n,m)+\underset{l=n}{\overset{\pm \mathrm{}}{}}K_\pm (n,l)\stackrel{~}{F}^\pm (l,m)=\frac{\delta (n,m)}{K_\pm (n,n)}\underset{j=1}{\overset{q}{}}\gamma _{\pm ,j}\widehat{\psi }_\pm (\rho _j,n)\widehat{\psi }_{q,\pm }(\rho _j,m)$$
and we have thus proved the following result.
###### Theorem 7.1.
The kernel $`K_\pm (n,m)`$ of the transformation operator satisfies the Gelโfand-Levitan-Marchenko equation
(7.6)
$$K_\pm (n,m)+\underset{l=n}{\overset{\pm \mathrm{}}{}}K_\pm (n,l)F^\pm (l,m)=\frac{\delta (n,m)}{K_\pm (n,n)},\pm m\pm n,$$
where
(7.7)
$$F^\pm (l,m)=\stackrel{~}{F}^\pm (l,m)+\underset{j=1}{\overset{q}{}}\gamma _{\pm ,j}\widehat{\psi }_{q,\pm }(\rho _j,l)\widehat{\psi }_{q,\pm }(\rho _j,m).$$
Defining the Gelโfand-Levitan-Marchenko operator
(7.8)
$$_n^\pm f(j)=\underset{l=0}{\overset{\mathrm{}}{}}F^\pm (n\pm l,n\pm j)f(l),f\mathrm{}^2(_0,),$$
yields that the Gelโfand-Levitan-Marchenko equation is equal to
(7.9)
$$(1+_n^\pm )K_\pm (n,n\pm .)=(K_\pm (n,n))^1\delta _0.$$
Our next aim is to study the Gelโfand-Levitan-Marchenko operator $`_n^\pm `$ in more detail. The structure of the Gelโfand-Levitan-Marchenko equation suggests that the estimate (5.5) for $`K_\pm (n,m)`$ should imply a similar estimate for $`F^\pm (n,m)`$.
###### Lemma 7.2.
(7.10)
$$|F^\pm (n,m)|C\underset{j=[\frac{n+m}{2}]\pm 1}{\overset{\pm \mathrm{}}{}}\left(|a(j)a_q(j)|+|b(j)b_q(j)|\right),$$
where the constant $`C`$ is of the same nature as in (5.5).
###### Proof.
We abbreviate the estimate (5.5) for $`K_+(n,m)`$ by
(7.11)
$$|K_+(n,m)|CC_+(n+m),$$
where
$$C_+(n+m)=\underset{j=[\frac{n+m}{2}]+1}{\overset{\mathrm{}}{}}c(j),c(j)=|a(j)a_q(j)|+|b(j)b_q(j)|.$$
Note that $`C_+(n+1)C_+(n)`$. Moreover, $`C_+(n)\mathrm{}_+^1()`$ since the summation by parts formula (e.g. , (1.18))
(7.12)
$$\underset{m=n}{\overset{N}{}}g(m)(f(m+1)f(m))=g(N)f(N+1)g(n1)f(n)+\underset{m=n}{\overset{N}{}}(g(m1)g(m))f(m)$$
implies for $`g(m)=m`$, $`f(m)=C_+(m)`$ that
(7.13)
$$\underset{m=n}{\overset{\mathrm{}}{}}mc(m)=(n1)C_+(n)+\underset{m=n}{\overset{\mathrm{}}{}}C_+(m),$$
where we used $`lim_n\mathrm{}nC_+(n+1)lim_n\mathrm{}_{m=n}^{\mathrm{}}mc(m)=0`$. Solving the GLM-equation (7.6) for $`F^+(n,m)`$, $`m>n`$, we obtain
$`|F^+(n,m)|`$ $``$ $`{\displaystyle \frac{1}{K_+(n,n)}}\left(|K_+(n,m)|+{\displaystyle \underset{l=n+1}{\overset{\mathrm{}}{}}}\left|K_+(n,l)F^+(l,m)\right|\right)`$
$``$ $`C_1(n)\left(C_+(n+m)+{\displaystyle \underset{l=n+1}{\overset{\mathrm{}}{}}}C_+(n+l)\left|F^+(l,m)\right|\right),`$
where $`C_1(n)=C|K_+(n,n)|^1C`$ for $`n\mathrm{}`$ (see (5.28)). For $`n`$ large enough, i.e. $`C_1(n)C_+(2n)<1`$, we apply the discrete Gronwall-type inequality , Lemma 10.8,
(7.14) $`|F^+(n,m)|`$ $``$ $`C_1(n)\left(C_+(n+m)+{\displaystyle \underset{l=n+1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{C_1(l)C_+(l+m)C_+(n+l)}{_{k=n+1}^l(1C_1(k)C_+(n+k))}}\right)`$
$``$ $`C_1(n)C_+(n+m)\left(1+{\displaystyle \underset{l=n+1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{C_1(k)C_+(n+l)}{_{k=n+1}^l(1C_1(n)C_+(n+k))}}\right),`$
which finishes the proof. โ
Furthermore,
###### Lemma 7.3.
Let $`F^\pm (n,m)`$ be solutions of the Gelโfand-Levitan-Marchenko equation. Then
(7.15) $`{\displaystyle \underset{n=n_0}{\overset{\pm \mathrm{}}{}}}|n|\left|F^\pm (n,n)F^\pm (n\pm 1,n\pm 1)\right|`$ $`<`$ $`\mathrm{},`$
(7.16) $`{\displaystyle \underset{n=n_0}{\overset{\pm \mathrm{}}{}}}|n|\left|a_q(n)F^\pm (n,n+1)a_q(n1)F^\pm (n1,n)\right|`$ $`<`$ $`\mathrm{}.`$
###### Proof.
We first prove (7.16) for $`F^+`$. Lemma 5.3 implies
(7.17)
$$b(n)b_q(n)=a_q(n)\kappa _{+,1}(n)a_q(n1)\kappa _{+,1}(n1),$$
where
(7.18)
$$\kappa _{+,j}(n):=\kappa _+(n,n+j):=\frac{K_+(n,n+j)}{K_+(n,n)}.$$
Abbreviate $`F_j^+(n):=F^+(n+j,n)`$. With this notation, the GLM-equation (7.6) reads
(7.19)
$$\kappa _{+,l}(n)+F_l^+(n)+\underset{j=1}{\overset{\mathrm{}}{}}\kappa _{+,j}(n)F_{jl}^+(n+l)=\frac{\delta (l,0)}{K_+(n,n)^2},l0.$$
Insert the GLM-equation for $`F^+(n,n+1)`$, $`F^+(n1,n)`$ (recall $`F^+(n,m)=F^+(m,n)`$)
(7.20) $`a_q(n)F_1^+(n)a_q(n1)F_1^+(n1)`$
$`=`$ $`a_q(n)\kappa _{+,1}(n)+a_q(n1)\kappa _{+,1}(n1)`$
$`{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}\left(a_q(n)\kappa _{+,j}(n)F_{j1}^+(n+1)a_q(n1)\kappa _{+,j}(n1)F_{j1}^+(n)\right).`$
Since $`a_q(n)\kappa _{+,1}(n)+a_q(n1)\kappa _{+,1}(n1)=b_q(n)b(n)`$ the only interesting part is the sum. For $`N`$, $`J<\mathrm{}`$,
(7.21) $`{\displaystyle \underset{n=n_0}{\overset{N}{}}}n{\displaystyle \underset{j=1}{\overset{J}{}}}\left(a_q(n)\kappa _{+,j}(n)F_{j1}^+(n+1)a_q(n1)\kappa _{+,j}(n1)F_{j1}^+(n)\right)`$
$`=`$ $`{\displaystyle \underset{j=1}{\overset{J}{}}}{\displaystyle \underset{n=n_0}{\overset{N}{}}}n\left(a_q(n)\kappa _{+,j}(n)F_{j1}^+(n+1)a_q(n1)\kappa _{+,j}(n1)F_{j1}^+(n)\right)`$
$`=`$ $`{\displaystyle \underset{j=1}{\overset{J}{}}}(Na_q(N)\kappa _{+,j}(N)F_{j1}^+(N+1)(n_01)a_q(n_01)\kappa _{+,j}(n_01)F_{j1}^+(n_0)`$
$`+{\displaystyle \underset{n=n_0}{\overset{N}{}}}(1)a_q(n1)\kappa _{+,j}(n1)F_{j1}^+(n)),`$
where we used the summation by parts. Estimates (7.11), (7.14) imply for the first summand
$`\left|{\displaystyle \underset{j=1}{\overset{J}{}}}Na_q(N)\kappa _{+,j}(N)F_{j1}^+(N+1)\right|`$ $``$ $`{\displaystyle \underset{j=1}{\overset{J}{}}}|N|a_q(N)\stackrel{~}{C}C_+(2N+j)C_+(2N+j+1)`$
$``$ $`|N|a_q(N)\widehat{C}C_+(2N+1),`$
which holds uniformly in $`J`$, and (compare (7.13))
(7.22)
$$\underset{N\mathrm{}}{lim}Na_q(N)\widehat{C}C_+(2N+1)=0.$$
Moreover,
$`\underset{N,J\mathrm{}}{lim}\left|{\displaystyle \underset{j=1}{\overset{J}{}}}{\displaystyle \underset{n=n_0}{\overset{N}{}}}a_q(n1)\kappa _{+,j}(n1)F_{j1}^+(n)\right|`$
$``$ $`\underset{N,J\mathrm{}}{lim}{\displaystyle \underset{j=1}{\overset{J}{}}}{\displaystyle \underset{n=n_0}{\overset{N}{}}}\left|a_q(n1)\kappa _{+,j}(n1)F_{j1}^+(n)\right|`$
$``$ $`{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=n_0}{\overset{\mathrm{}}{}}}a_q(n1)\stackrel{~}{C}C_+(2n+j)C_+(2n+j+1)<\mathrm{}.`$
Therefore $`|n||a_q(n)F^+(n,n+1)a_q(n1)F^+(n1,n)|\mathrm{}_+^1()`$ as desired. To apply Lemma 5.3 for $`F^{}`$ use the symmetry property $`F^{}(n,m)=F^{}(m,n)`$. For (7.15), inserting the GLM-equation yields
$`F^+(n,n)F^+(n+1,n+1)=K_+^2(n,n)K_+^2(n+1,n+1)`$
$`+{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}\left(\kappa _{+,j}(n+1)F_j^+(n+1)\kappa _{+,j}(n)F_j^+(n)\right).`$
By (5.28),
(7.23) $`\left|K_+^2(n,n)K_+^2(n+1,n+1)\right|`$ $``$ $`{\displaystyle \frac{|a(n)+a_q(n)|}{a(n)^2}}{\displaystyle \underset{j=n+1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{a(j)^2}{a_q(j)^2}}|a(n)a_q(n)|`$
$``$ $`C|a(n)a_q(n)|,`$
and the same considerations as above imply (7.15). โ
###### Remark 7.4.
The Gelโfand-Levitan-Marchenko equation is symmetric in $`K_\pm (n,m)`$ and $`F^\pm (n,m)`$, therefore we can invert the analysis done in Lemma 7.3 and obtain estimates for $`K_\pm (n,m)`$ starting with an analogue of estimate (7.10) for $`F^\pm (n,m)`$ and the estimates (7.15), (7.16) (c.f. Lemma 8.1).
###### Theorem 7.5.
For $`n`$, the Gelโfand-Levitan-Marchenko operator $`_n^\pm :\mathrm{}^2\mathrm{}^2`$ is Hilbert-Schmidt. Moreover, $`1+_n^\pm `$ is positive and hence invertible.
In particular, the Gelโfand-Levitan-Marchenko equation (7.9) has a unique solution and $`S_+(H)`$ or $`S_{}(H)`$ uniquely determine $`H`$.
###### Proof.
That $`_n^\pm `$ is Hilbert-Schmidt is a straight-forward consequence of our estimate Lemma 7.2.
Let $`f\mathrm{}^2(_0)`$ be real (which is no restriction since $`F^+(n,l)`$ is real and the real and imaginary part of (7.24) could be treated separately) and abbreviate $`f_n(w)=_{j=0}^{\mathrm{}}f(j)\psi _{q,+}(w,n+j)`$. Then
$`{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}f(j)_n^+f(j)={\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}f(j){\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}F^+(n+j,n+l)f(l)`$
$`={\displaystyle \frac{1}{2\pi i}}{\displaystyle _{|w|=1}}R_+(w){\displaystyle \underset{j,l=0}{\overset{\mathrm{}}{}}}f(j)\psi _{q,+}(w,n+j)\psi _{q,+}(w,n+l)f(l)d\omega (w)`$
$`+{\displaystyle \underset{k=1}{\overset{q}{}}}{\displaystyle \underset{j,l=0}{\overset{\mathrm{}}{}}}f(j)\gamma _{+,k}\widehat{\psi }_{q,+}(\rho _k,n+j)\widehat{\psi }_{q,+}(\rho _k,n+l)f(l)`$
$`={\displaystyle \frac{1}{2\pi i}}{\displaystyle _{|w|=1}}R_+(w)f_n(\overline{w})f_n(w)๐\omega (w)+{\displaystyle \underset{k=1}{\overset{q}{}}}\gamma _{+,k}|\widehat{f}_n(\rho _k)|^2`$
(7.24) $`={\displaystyle \frac{1}{2\pi i}}{\displaystyle _{|w|=1}}\stackrel{~}{R}_+(w)|f_n(w)|^2๐\omega (w)+{\displaystyle \underset{k=1}{\overset{q}{}}}\gamma _{+,k}|\widehat{f}_n(\rho _k)|^2,`$
where $`\stackrel{~}{R}_+(w)=R_+(w)f_n(w)\left(\overline{f_n(w)}\right)^1`$ with $`|\stackrel{~}{R}_+(w)|=|R_+(w)|`$ and $`\widehat{f}_n(w)=_{j=0}^{\mathrm{}}f(j)\widehat{\psi }_{q,+}(w,n+j)`$. The integral over the imaginary part vanishes since $`\overline{\stackrel{~}{R}_+(w)}=\stackrel{~}{R}_+(\overline{w})`$ and we replace the real part by
$$\text{Re}(\stackrel{~}{R}_+(w))=\frac{1}{2}\left(|1+\stackrel{~}{R}_+(w)|^21|\stackrel{~}{R}_+(w)|^2\right)=\frac{1}{2}\left(|1+\stackrel{~}{R}_+(w)|^2+|T(w)|^2\right)1,$$
(recall $`|\stackrel{~}{R}_+(w)|^2+|T(w)|^2=1`$). This yields using $`|f(j)|^2=\frac{1}{2\pi i}_{|w|=1}|f_n(w)|^2๐\omega `$
(7.25) $`{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}\overline{f(j)}(1+_n^+)f(j)`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{q}{}}}\gamma _{+,k}|\widehat{f}_n(\rho _k)|^2`$
$`+{\displaystyle \frac{1}{4\pi i}}{\displaystyle _{|w|=1}}\left(|1+\stackrel{~}{R}_+(w)|^2+|T(w)|^2\right)|f_n(w)|^2๐\omega (w)`$
which establishes $`1+_n^+0`$. According to Lemma 6.2, $`|T(w)|^2>0`$ a.e., therefore $`1`$ is not an eigenvalue and $`1+_n^+ฯต_n`$ for some $`ฯต_n>0`$. โ
To finish the direct scattering step for the Jacobi operator $`H`$ with asymptotically quasi-periodic coefficients we summarize the properties of the scattering data $`S_\pm (H)`$.
###### Hypothesis H. 7.6.
The scattering data
(7.26)
$$S_\pm (H)=\{R_\pm (w),|w|=1;(\rho _j,\gamma _{\pm ,j}),1jq\}$$
satisfy the following conditions.
(i). The reflection coefficients $`R_\pm (w)`$ are continuous except possibly at $`w_l=w(E_l)`$ and fulfill
(7.27)
$$\overline{R_\pm (w)}=R_\pm (\overline{w}).$$
Moreover, $`|R_\pm (w)|<1`$ for $`ww_l`$ and
(7.28)
$$1|R_\pm (w)|^2C\underset{l=0}{\overset{2g+1}{}}|ww_l|^2.$$
The Fourier coefficients
(7.29)
$$\stackrel{~}{F}^\pm (l,m)=\frac{1}{2\pi \mathrm{i}}_{|w|=1}R_\pm (w)\psi _{q,\pm }(w,l)\psi _{q,\pm }(w,m)๐\omega (w)$$
satisfy
$`|\stackrel{~}{F}^\pm (n,m)|{\displaystyle \underset{j=n+m}{\overset{\pm \mathrm{}}{}}}q(j),q(j)0,|j|q(j)\mathrm{}^1(),`$
$`{\displaystyle \underset{n=n_0}{\overset{\pm \mathrm{}}{}}}|n|\left|\stackrel{~}{F}^\pm (n,n)\stackrel{~}{F}^\pm (n\pm 1,n\pm 1)\right|<\mathrm{},`$
$`{\displaystyle \underset{n=n_0}{\overset{\pm \mathrm{}}{}}}|n|\left|a_q(n)\stackrel{~}{F}^\pm (n,n+1)a_q(n1)\stackrel{~}{F}^\pm (n1,n)\right|<\mathrm{}.`$
(ii). The values $`\rho _j\backslash \sigma (H_q)`$, $`1jq`$, are distinct and the norming constants $`\gamma _{\pm ,j}`$, $`1jq`$, are positive.
(iii). $`T(w)`$ defined via equation (6.25) extends to a single valued function on $`๐`$ (i.e., it has equal values on the corresponding slits).
(iv). Transmission and reflection coefficients satisfy and satisfies
(7.30)
$$\begin{array}{cc}\underset{ww_l}{lim}(ww_l)\frac{R_\pm (w)+1}{T(w)}=0,\hfill & w_lw(\mu _j)\hfill \\ \underset{ww_l}{lim}(ww_l)\frac{R_\pm (w)1}{T(w)}=0\hfill & w_l=w(\mu _j)\hfill \end{array}.$$
and the consistency conditions
$$\frac{R_{}(w)}{R_+(\overline{w})}=\frac{T(w)}{T(\overline{w})},\gamma _{+,j}\gamma _{,j}=\frac{\left(\mathrm{Res}_{\rho _j}T(\lambda )\right)^2}{_{l=0}^{2g+1}(\rho _jE_l)}.$$
###### Remark 7.7.
Note that (7.28) implies that $`\mathrm{ln}(1|R_\pm (w)|^2)`$ is integrable and ensures that (6.25) is well-defined, at least as a multi valued function. Condition (iii), which is void in the constant background case, shows that the the reflection coefficient and eigenvalues cannot be chosen independent of each other.
## 8. Inverse scattering theory
In this section we want to invert the process of scattering theory, that is, we want to reconstruct the operator $`H`$ from a given set $`S_\pm `$ and a given quasi-periodic Jacobi operator $`H_q`$.
If $`S_\pm `$ (satisfying H.7.6 (i)โ(ii)) and $`H_q`$ are known, we can construct $`F^\pm (l,m)`$ via formula (7.7) and thus derive the Gelโfand-Levitan-Marchenko equation, which has a unique solution by Theorem 7.5. This solution
$`K_\pm (n,n)`$ $`=`$ $`\delta _0,(1+_n^\pm )^1\delta _0^{1/2}`$
(8.1) $`K_\pm (n,n\pm j)`$ $`=`$ $`{\displaystyle \frac{1}{K_\pm (n,n)}}\delta _j,(1+_n^\pm )^1\delta _0`$
is the kernel of the transformation operator. Since $`1+_n^\pm `$ is positive, $`K_\pm (n,n)`$ is positive and we can set in accordance with Lemma 5.3
(8.2) $`a_+(n)`$ $`=`$ $`a_q(n){\displaystyle \frac{K_+(n+1,n+1)}{K_+(n,n)}},`$
$`a_{}(n)`$ $`=`$ $`a_q(n){\displaystyle \frac{K_{}(n,n)}{K_{}(n+1,n+1)}},`$
$`b_+(n)`$ $`=`$ $`b_q(n)+a_q(n){\displaystyle \frac{K_+(n,n+1)}{K_+(n,n)}}a_q(n1){\displaystyle \frac{K_+(n1,n)}{K_+(n1,n1)}},`$
$`b_{}(n)`$ $`=`$ $`b_q(n)+a_q(n1){\displaystyle \frac{K_{}(n,n1)}{K_{}(n,n)}}a_q(n){\displaystyle \frac{K_{}(n+1,n)}{K_{}(n+1,n+1)}}.`$
Let $`H_+`$, $`H_{}`$ be the associated Jacobi operators.
###### Lemma 8.1.
Suppose a given set $`S_\pm `$ satisfies H.7.6 (i)โ(ii). Then the sequences defined in (8.2) satisfy $`n|a_\pm (n)a_q(n)|`$, $`n|b_\pm (n)b_q(n)|\mathrm{}_\pm ^1()`$.
Moreover, $`\psi _\pm (\lambda ,n)=_{m=n}^\pm \mathrm{}K_\pm (n,m)\psi _{q,\pm }(\lambda ,m)`$, where $`K_\pm (n,m)`$ is the solution of the Gelโfand-Levitan-Marchenko equation, satisfies $`\tau _\pm \psi _\pm =\lambda \psi _\pm `$.
###### Proof.
We only prove the statements for the โ+โ case. Define $`F^+(n,m)`$ by (c.f. (7.7))
$$F^+(l,m)=\stackrel{~}{F}^+(l,m)+\underset{j=1}{\overset{q}{}}\gamma _{+,j}\widehat{\psi }_{q,+}(\rho _j,l)\widehat{\psi }_{q,+}(\rho _j,m).$$
Hypothesis H.7.6 (i) implies
(8.3) $`|F^+(n,m)|C{\displaystyle \underset{j=n+m}{\overset{\mathrm{}}{}}}q(j)=:C_+(n+m),`$
(8.4) $`{\displaystyle \underset{n=n_0}{\overset{\mathrm{}}{}}}|n|\left|F^+(n,n)F^+(n+1,n+1)\right|<\mathrm{},`$
(8.5) $`{\displaystyle \underset{n=n_0}{\overset{\mathrm{}}{}}}|n|\left|a_q(n)F^+(n,n+1)a_q(n1)F^+(n1,n)\right|<\mathrm{},`$
since $`\widehat{\psi }_{q,+}(\rho _j,n)`$ decay exponentially as $`n\mathrm{}`$ and $`_j\gamma _{+,j}\widehat{\psi }_{q,+}(\rho _j,.)\widehat{\psi }_{q,+}(\rho _j,.)`$ form a telescopic sum. Note that $`C_+(n+1)<C_+(n)`$.
Set $`\kappa _+(n,m):=K_+(n,m)K_+(n,n)^1`$. Then as in the proof of Lemma 7.2 we obtain
(8.6)
$$|\kappa _+(n,m)|C_+(n+m)(1+O(1)).$$
Now we have all estimates at our disposal to prove $`n|b_+(n)b_q(n)|\mathrm{}^1()`$. By definition (c.f. (8.2)),
(8.7)
$$b_+(n)b_q(n)=a_q(n)\kappa _+(n,n+1)a_q(n1)\kappa _+(n1,n).$$
We insert the GLM-equation for $`\kappa _+(n,n+1)`$, $`\kappa _+(n1,n)`$ and use estimate (8.5), the summation by parts formula, and estimates (8.3), (8.6) in the same way as in Lemma 7.3. Similarly using (8.4) we see
(8.8)
$$\underset{n=n_0}{\overset{\mathrm{}}{}}|n|\left|\frac{1}{K_+^2(n,n)}\frac{1}{K_+^2(n+1,n+1)}\right|<\mathrm{}.$$
Equation (8.2) yields
$$\left|\frac{1}{K_+^2(n,n)}\frac{1}{K_+^2(n+1,n+1)}\right|=\frac{1}{a_q(n)^2}\left(\underset{j=n+1}{\overset{\mathrm{}}{}}\frac{a_+(j)^2}{a_q(j)^2}\right)|a_+(n)^2a_q(n)^2|.$$
The product converges and therefore $`|n||a_+(n)^2a_q(n)^2|\mathrm{}^1()`$.
Next we consider $`\psi _+(\lambda ,n)`$. Abbreviate
(8.9) $`(\mathrm{\Delta }K_+)(n,m)=a_q(n1)\kappa _+(n1,m)+a_+^2(n)a_q^1(n)\kappa _+(n+1,m)`$
$`a_q(m1)\kappa _+(n,m1)a_q(m)\kappa _+(n,m+1)+(b_+(n)b_q(m))\kappa _+(n,m).`$
$`\mathrm{\Delta }K_+=0`$ is equivalent to the operator equality $`H_+K_+=K_+H_q`$, which in turn implies that $`\psi _+(\lambda ,n)`$ satisfies $`H_+\psi _+=\lambda \psi _+`$
(8.10)
$$H_+\psi _+=H_+K_+\psi _{q,+}=K_+H_q\psi _{q,+}=K_+\lambda \psi _{q,+}=\lambda K_+\psi _{q,+}=\lambda \psi _+.$$
To show that $`\mathrm{\Delta }K_+=0`$ we insert the GLM-equation into (8.9) and obtain
(8.11)
$$(\mathrm{\Delta }K_+)(n,m)+\underset{l=n+1}{\overset{\mathrm{}}{}}(\mathrm{\Delta }K_+)(n,l)F^+(l,m)=0,m>n+1.$$
In the calculations we used
$`a_q(n1)F^+(n1,m)+b_q(n)F^+(n,m)+a_q(n)F^+(n+1,m)=`$
$`a_q(m1)F^+(n,m1)+b_q(m)F^+(n,m)+a_q(m)F^+(n,m+1)`$
which follows from (7.7). By Theorem 7.5 equation (8.11) has only the trivial solution $`\mathrm{\Delta }K_+=0`$ and hence the proof is complete. โ
Now we can prove the main result of this section.
###### Theorem 8.2.
Hypothesis H.7.6 is necessary and sufficient for a sets $`S_\pm `$ to be the left/right scattering data of a unique Jacobi operator $`H`$ associated with sequences $`a`$, $`b`$ satisfying H.4.1.
###### Proof.
Necessity has been established in the previous section. By Lemma 8.1, we know existence of sequences $`a_\pm `$, $`b_\pm `$ and corresponding solutions $`\psi _\pm (w,n)`$ associated with $`S_+`$ (or $`S_{}`$). Hence it remains to establish $`a_+(n)=a_{}(n)`$ and $`b_+(n)=b_{}(n)`$.
Consider the following part of the GLM-equation
(8.12)
$$\mathrm{\Phi }_+(n,.):=\underset{l=n}{\overset{\mathrm{}}{}}K_+(n,l)\stackrel{~}{F}^+(l,.)\mathrm{}_+^1().$$
Then by use of (7.2) and Lemma 3.6,
(8.13) $`{\displaystyle \underset{m}{}}\mathrm{\Phi }_+(n,m)\psi _{q,}(w,m)={\displaystyle \underset{m}{}}\left({\displaystyle \underset{l=n}{\overset{\mathrm{}}{}}}K_+(n,l)\stackrel{~}{F}^+(l,m)\right)\psi _{q,}(w,m)`$
$`=`$ $`{\displaystyle \underset{m}{}}\left({\displaystyle \frac{1}{2\pi \mathrm{i}}}{\displaystyle _{|w|=1}}R_+(w)\psi _+(w,n)\psi _{q,+}(w,m)๐\omega (w)\right)\psi _{q,}(w,m)`$
$`=`$ $`{\displaystyle \underset{m}{}}\psi _{q,}(w,m),R_+(w)\psi _+(w,n)\psi _{q,}(w,m)`$
$`=`$ $`R_+(w)\psi _+(w,n).`$
On the other hand, inserting the GLM-equation yields for $`|w|=1`$
(8.14) $`{\displaystyle \underset{m}{}}\mathrm{\Phi }_+(n,m)\psi _{q,}(w,m)=`$
$`=`$ $`{\displaystyle \underset{m=\mathrm{}}{\overset{n1}{}}}\mathrm{\Phi }_+(n,m)\psi _{q,}(w,m)+{\displaystyle \underset{m=n}{\overset{\mathrm{}}{}}}[\delta (n,m)K_+^1(n,n)K_+(n,m)`$
$`{\displaystyle \underset{l=n}{\overset{\mathrm{}}{}}}K_+(n,l){\displaystyle \underset{j=1}{\overset{q}{}}}\gamma _{+,j}\widehat{\psi }_{q,+}(\rho _j,l)\widehat{\psi }_{q,+}(\rho _j,m)]\psi _{q,}(w,m)`$
$`=`$ $`{\displaystyle \underset{m=\mathrm{}}{\overset{n1}{}}}\mathrm{\Phi }_+(n,m)\psi _{q,}(w,m)+\psi _{q,}(w,n)K_+^1(n,n)\overline{\psi _+(w,n)}`$
$`{\displaystyle \underset{j=1}{\overset{q}{}}}\gamma _{+,j}\widehat{\psi }_+(\rho _j,n){\displaystyle \underset{m=n}{\overset{\mathrm{}}{}}}\widehat{\psi }_{q,+}(\rho _j,m)\psi _{q,}(w,m),`$
(recall the definition of $`\widehat{\psi }_{q,\pm }`$ from (6.8)) and therefore
(8.15)
$$T(w)h_{}(w,n)=\overline{\psi _+(w,n)}+R_+(w)\psi _+(w,n),|w|=1,$$
where
(8.16) $`h_{}(w,n)`$ $`=`$ $`{\displaystyle \frac{\psi _{q,}(w,n)}{T(w)}}({\displaystyle \frac{1}{K_+(n,n)}}+{\displaystyle \underset{m=\mathrm{}}{\overset{n1}{}}}\mathrm{\Phi }_+(n,m){\displaystyle \frac{\psi _{q,}(w,m)}{\psi _{q,}(w,n)}}`$
$`+{\displaystyle \underset{j=1}{\overset{q}{}}}\gamma _{+,j}\widehat{\psi }_+(\rho _j,n){\displaystyle \frac{W_{n1}(\widehat{\psi }_{q,+}(\rho _j),\psi _{q,}(w))}{\psi _{q,}(w,n)(\lambda (w)\rho _j)}}),`$
since Greenโs formula (, eq. (1.20)) implies for $`\lambda \sigma (H_q)`$
$$(\lambda \rho _j)\underset{m=n}{\overset{\mathrm{}}{}}\widehat{\psi }_{q,+}(\rho _j,m)\psi _{q,}(\lambda ,m)=W_{n1}(\widehat{\psi }_{q,+}(\rho _j),\psi _{q,}(\lambda )).$$
Similarly, we obtain
(8.17) $`h_+(w,n)`$ $`=`$ $`{\displaystyle \frac{\psi _{q,+}(w,n)}{T(w)}}({\displaystyle \frac{1}{K_{}(n,n)}}+{\displaystyle \underset{m=n+1}{\overset{\mathrm{}}{}}}\mathrm{\Phi }_{}(n,m){\displaystyle \frac{\psi _{q,+}(w,m)}{\psi _{q,+}(w,n)}}`$
$`{\displaystyle \underset{j=1}{\overset{q}{}}}\gamma _{,j}\widehat{\psi }_{}(\rho _j,n){\displaystyle \frac{W_n(\widehat{\psi }_{q,}(\rho _j),\psi _{q,+}(w))}{\psi _{q,+}(w,n)(\lambda (w)\rho _j)}})`$
with
$$\mathrm{\Phi }_{}(n,m)=\underset{l=\mathrm{}}{\overset{n}{}}K_{}(n,l)\stackrel{~}{F}^{}(l,m).$$
For $`n`$, $`|w|=1`$, we see that $`h_{}(w^1,n)=\overline{h_{}(w,n)}`$, since $`K_\pm (n,m)`$ and $`\mathrm{\Phi }_\pm (n,m)`$ are real. The functions $`h_{}(w,n)`$ are continuous for $`|w|=1`$, $`ww(E_j)`$, since $`T^1(w)`$ is continuous on this set by the Poisson-Jensen formula (6.25) ($`|R_\pm (w)|<1`$ for $`ww(E_j)`$ by H.7.6 (i)) and $`\psi _{q,}(w,m)`$ are continuous on $`W\backslash \{w(\mu _k)\}`$. The functions $`h_{}(w,n)`$ have a meromorphic continuation to $`๐\backslash \{0\}`$ with the only possible poles at $`w(\rho _j)`$ and $`w(\mu _j)`$. At $`w(\rho _j)`$ there are no poles, due to the zeros of $`T^1(w)`$ at $`w(\rho _j)`$. For $`w=w(\mu _j)`$ we have the same type of singularity as $`\psi _{q,\pm }`$. In summary, $`h_\pm (w,n)`$ have simple poles at $`w(\mu _j)`$ and are continuous at the boundary except possibly at $`w(E_j)`$.
To study the behavior of $`h_\pm (w,n)`$ as $`w0`$, we recall $`z^1=w/\stackrel{~}{a}(1+O(w))`$. Then
$`\left({\displaystyle \frac{w}{\stackrel{~}{a}}}+O(w^2)\right)W_{n1}(\widehat{\psi }_{q,+}(\rho _j),\psi _{q,}(w))`$
$`=`$ $`{\displaystyle \frac{(1)^n\stackrel{~}{a}^{n1}}{_{j=0}^{n2}a_q(j)}}w^{n+1}(\widehat{\psi }_{q,+}(\rho _j,n1)+O(w)),`$
$`\left({\displaystyle \frac{w}{\stackrel{~}{a}}}+O(w^2)\right)W_n(\widehat{\psi }_{q,}(\rho _j),\psi _{q,+}(w))`$
$`=`$ $`{\displaystyle \frac{(1)^n_{j=0}^na_q(j)}{\stackrel{~}{a}^{n+1}}}w^{n+1}(\widehat{\psi }_{q,}(\rho _j,n+1)+O(w)),`$
and property (B4) implies
(8.18)
$$\underset{m=n1}{\overset{\mathrm{}}{}}\mathrm{\Phi }_\pm (n,m)\psi _{q,}(w,m)\psi _{q,}^1(w,n)=O(w),w0.$$
We conclude that
(8.19)
$$\underset{w0}{lim}h_{}(w,n)\psi _{q,\pm }(w,n)=\frac{1}{T(0)K_\pm (n,n)}.$$
H.7.6 (iv) and (6.1) imply the following behavior of $`\widehat{h}_{}(\lambda ,n)`$ as $`\lambda \rho _j`$
(8.20) $`\underset{\lambda \rho _j}{lim}\widehat{h}_{}(\lambda ,n)`$ $`=`$ $`\pm \gamma _{\pm ,j}\widehat{\psi }_\pm (\rho _j,n)\underset{\lambda \rho _j}{lim}{\displaystyle \frac{W_{n1}(\widehat{\psi }_{q,\pm }(\rho _j),\widehat{\psi }_{q,}(\lambda ))}{(\lambda \rho _j)T(\lambda )}}`$
$`=`$ $`\gamma _{\pm ,j}\widehat{\psi }_\pm (\rho _j,n)\left(\mathrm{Res}_{\rho _j}T(\lambda )\right)^1{\displaystyle \underset{l=0}{\overset{2g+1}{}}}\sqrt{\rho _jE_l},`$
where $`\widehat{h}_\pm `$ are defined as in (6.8).
By virtue of the consistency condition $`T(w)\overline{R_+(w)}=\overline{T(w)}R_{}(w)`$ we obtain
$`\overline{h_\pm (w,n)}+R_\pm (w)h_\pm (w,n)=`$
$`=`$ $`{\displaystyle \frac{1}{\overline{T(w)}}}\left(\psi _{}(w,n)+\overline{R_{}(w)}\overline{\psi _{}(w,n)}\right)+{\displaystyle \frac{R_\pm (w)}{T(w)}}\left(\overline{\psi _{}(w,n)}+R_{}(w)\psi _{}(w,n)\right)`$
$`=`$ $`\psi _{}(w,n)\left({\displaystyle \frac{1}{\overline{T(w)}}}+{\displaystyle \frac{R_\pm (w)R_{}(w)}{T(w)}}\right)+\overline{\psi _{}(w,n)}\left({\displaystyle \frac{\overline{R_{}(w)}}{\overline{T(w)}}}+{\displaystyle \frac{R_\pm (w)}{T(w)}}\right)`$
$`=`$ $`\psi _{}(w,n)T(w),|w|=1.`$
If we eliminate $`R_{}(w)`$ from the last equation and (8.15) we see
$`T(w)R_{2g+2}^{1/2}(w)\left(\widehat{\psi }_+(w,n)\widehat{\psi }_{}(w,n)\widehat{h}_+(w,n)\widehat{h}_{}(w,n)\right)`$ $`=`$
(8.21) $`{\displaystyle \frac{_j(\lambda (w)\mu _j)}{R_{2g+2}^{1/2}(w)}}\left(\overline{h_\pm (w,n)}\psi _\pm (w,n)\overline{\psi _\pm (w,n)}h_\pm (w,n)\right)`$ $`=:`$ $`G(w,n),`$
for $`|w|=1`$. Observe that $`G(\overline{w},n)=\overline{G(w,n)}=G(w,n)`$, $`|w|=1`$, since $`\overline{\widehat{h}_\pm }\widehat{\psi }_\pm \overline{\widehat{\psi }_\pm }\widehat{h}_\pm `$ and $`R_{2g+2}^{1/2}(w)`$ are odd functions for $`|w|=1`$. The function $`G(w,n)`$ can be continued analytically on $`๐`$ since the difference $`\widehat{\psi }_+\widehat{\psi }_{}\widehat{h}_+\widehat{h}_{}`$ vanishes at the poles $`w(\rho _j)`$ of $`T(w)`$ by (8.20). Note that the product $`\widehat{\psi }_+\widehat{\psi }_{}`$ and hence also $`\widehat{h}_+\widehat{h}_{}`$ do not have poles at $`w(\mu _j)`$. Moreover, since $`๐`$ is just the image of the upper sheet, we can extend it to a compact Riemann surface $`\stackrel{~}{๐}`$ by adding the image of the lower sheet. Now by $`G(\overline{w},n)=G(w,n)`$ we can extend $`G`$ to $`\stackrel{~}{๐}`$ by setting $`G(w,n)=G(w^1,n)`$ for $`|w|>1`$.
Now let us investigate the behavior at the band edges: If $`w_lw(\mu _j)`$, we obtain by (7.30), (8.15), and real-valuedness of $`\widehat{\psi }_\pm `$ at the band edges that
$`\underset{ww_l}{lim}R_{2g+2}^{1/2}(w){\displaystyle \underset{j}{}}(\lambda (w)\mu _j)h_{}(w,n)\overline{\psi _{}(w,n)}`$
$`=`$ $`\underset{ww_l}{lim}{\displaystyle \frac{R_{2g+2}^{1/2}_j(\lambda \mu _j)}{T}}\left(\overline{\psi _\pm }+R_\pm \psi _\pm \right)\overline{\psi _{}}`$
$`=`$ $`\underset{ww_l}{lim}{\displaystyle \frac{R_{2g+2}^{1/2}_j(\lambda \mu _j)}{T}}\left((R_\pm +1)\psi _\pm +\overline{\psi _\pm }\psi _\pm \right)\overline{\psi _{}}=0.`$
If $`w_l=w(\mu _j)`$, the same calculation shows that
$`\underset{ww_l}{lim}R_{2g+2}^{1/2}(w){\displaystyle \underset{j}{}}(\lambda (w)\mu _j)h_\pm (w,n)\overline{\psi _\pm (w,n)}`$
$`=(1)^{l+1}C_+(n)C_{}(n)\underset{ww_l}{lim}R_{2g+2}^{1/2}(w){\displaystyle \frac{R_\pm (w)1}{T(w)}}=0`$
by (7.30), where we used $`\psi _\pm (w,n)=\mathrm{i}^lC_\pm (n)(\lambda (w)\mu _j)^{1/2}+O(1)`$.
Consequently $`R_{2g+2}(w)G(w,n)`$ is continuous at $`w=w_l`$ and vanishes at the band edges. Thus the singularities of $`R_{2g+2}^{1/2}(w)G(w,n)`$ at $`w_l`$ are removable. Furthermore, $`R_{2g+2}^{1/2}(w)G(w,n)`$ is purely imaginary for $`|w|=1`$ and real on the slits and hence must vanish at $`w_l`$ by continuity. So the singularities of $`G(w,n)`$ at $`w_l`$ are removable as well. Thus $`G`$ is holomorphic on all of $`\stackrel{~}{๐}`$ and vanishes at $`w=0`$, that is, $`G(w,n)0`$ which implies (compare (B4))
$`\underset{w0}{lim}\left(\psi _+(w,n)\psi _{}(w,n)h_+(w,n)h_{}(w,n)\right)`$
$`=`$ $`K_+(n,n)K_{}(n,n)(T(0)^2K_+(n,n)K_{}(n,n))^1=0.`$
Using (8.2) we finally obtain from $`T(0)^2=\left(K_+(n,n)K_{}(n,n)\right)^2`$ that
(8.22)
$$a_+(n)=a_{}(n)a(n),n.$$
It remains to prove $`b_+(n)=b_{}(n)`$. Proceeding as for $`G(w,n)`$ we can show that
$`T(w)R_{2g+2}^{1/2}(w)\left(\widehat{\psi }_+(w,n)\widehat{\psi }_{}(w,n+1)\widehat{h}_+(w,n+1)\widehat{h}_{}(w,n)\right)=`$
(8.23) $`{\displaystyle \frac{_j(\lambda (w)\mu _j)}{R_{2g+2}^{1/2}(w)}}\left(\overline{h_+(w,n+1)}\psi _+(w,n)\overline{\psi _+(w,n)}h_+(w,n+1)\right)`$
is a constant equal to $`1/a(n)`$. Thus
(8.24) $`W(w,n)`$ $`:=`$ $`a(n)\left(\psi _+(w,n)\psi _{}(w,n+1)h_+(w,n+1)h_{}(w,n)\right)`$
$`=`$ $`{\displaystyle \frac{R_{2g+2}^{1/2}(w)}{T(w)_j(\lambda (w)\mu _j)}}.`$
and computing the asymptotics at $`w=0`$ (compare (4.3)) we see
(8.25)
$$0=W(w,n)W(w,n1)=A_+(0)A_{}(0)(b_+(n)b_{}(n))$$
and in particular $`b_+(n)=b_{}(n)b(n)`$.
Our operator $`H`$ has the correct norming constants since as in (6.12) it follows
(8.26)
$$\underset{n}{}\widehat{\psi }_+(\rho _j,n)\widehat{\psi }_{}(\rho _j,n)=\left(\mathrm{Res}_{\rho _j}T(\lambda )\right)^1\underset{l=0}{\overset{2g+1}{}}\sqrt{\rho _jE_l}$$
and by (8.20),
$$\underset{n}{}\widehat{\psi }_\pm (\rho _j,n)\widehat{\psi }_\pm (\rho _j,n)=\gamma _{\pm ,j}^1.$$
## Acknowledgments
I.E. thanks A. Boutet de Monvel for the kind hospitality of University Paris-7, where part of this work was done. G.T. thanks Peter Yuditskii for several helpful discussions and hints with respect to the literature. We thank Mark Losik for help with respect to literature.
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# The ideal glass transition of Hard Spheres
## I Introduction
The question whether a liquid of identical Hard Spheres undergoes a glass transition upon densification is still open RT96 ; RLSB98 ; Sp98 ; TdCFNC04 . If crystallization is avoided, one can access the metastable region of the phase diagram above the freezing packing fraction $`\phi _f=0.494`$, where $`\phi =\frac{N\pi D^3}{6V}`$, $`D`$ is the Hard Sphere diameter, $`N`$ is the number of particles and $`V`$ is the volume of the container. In this region the dynamics of the liquid becomes slower and slower on increasing the density. The particles are โcagedโ by their neighbors, and the dynamics separates into a fast rattling inside the cage and slow rearrangements of the cages. The typical time scale of these rearrangements increase very fast around $`\phi _g0.56`$ and many authors reported the observation of a glass transition at these values of density GS91 ; vMU93 .
If the radius of the cages is sufficiently small and if the typical time scale of cage rearrangements is sufficiently large, the system vibrates around configurations that are stable for a very large time and can be threated as metastable states. It is then natural to separate the total entropy of the liquid in a โvibrationalโ contribution, that accounts for the entropy related to the rattling of the particles around the metastable structure, and a โconfigurationalโ entropy that is the number of metastable states accessible to the liquid at the considered value of density SW84 ; DeB96 . For many simple potentials such as the LennardโJones CMPV99 ; SKT99 and for more realistic systems as well Ka48 ; An95 the extrapolation of the measured configurational entropy at higher density (or lower temperature) indicates that there exists a density, called Kauzmann density $`\phi _K`$, where the configurational entropy vanishes. The system freezes in the lowest free-energy states and no more rearrangements of the structure are possible. This transition is commonly called ideal glass transition or Kauzmann transition DeB96 ; CMPV99 ; SKT99 ; Ka48 ; An95 ; MP99 . Note that the Kauzmann density is expected to be larger than the experimental glass transition density, as at $`\phi _K`$ the relaxation time is expected to diverge so that the system freezes in a metastable state, on the experimental time scale, for a density $`\phi _g`$ smaller than $`\phi _K`$. The density $`\phi _g`$ where the real glass transition happens (weakly) depends on the experimentally accessible time scale. Few estimates of the configurational entropy for Hard Spheres are currently available Sp98 ; CFP98 ; Luca05 and indicate a value of $`\phi _K`$ in the range $`0.58รท0.62`$.
A related problem is the study of dense amorphous packings of Hard Spheres. Dense amorphous packings are relevant in the study of colloidal suspensions, granular matter, powders, etc. and have been widely studied in the literature Be83 ; SK69 ; Fi70 ; Be72 ; Ma74 ; Po79 ; Al98 ; SEGHL02 . The amorphous metastable configurations described above provide examples of such packings: when the system freezes in one of these states, if one is still able to increase the density in order to reduce the size of the cages to zero (for example by shaking the container SK69 ; Fi70 or making use of suitable computer algorithms Be72 ; Ma74 ; SEGHL02 ), a random close packed state is reached. The problem of which is the maximum value of density $`\phi _c`$ that can be reached applying this kind of procedures has been tackled using a lot of different techniques, usually finding values of $`\phi _c`$ in the range $`0.62รท0.67`$. Another interesting problem is to estimate the mean coordination number $`z`$, i.e. the mean number of contacts between a sphere and its neighbors, in the random close packed states. Many studies addressed this question usually finding values of $`z6`$.
Recently, the replica method MPV87 ; Mo95 ; MP99 has been successfully applied to the study of the ideal glass transition in simple liquids as the LennardโJones liquid. Reliable estimates of the configurational entropy, of the Kauzmann temperature and of the thermodynamic properties of the glass have been obtained from first principles in this way MP99 ; MP99b ; CMPV99 ; MP00 . However, for technical reasons this approach could not be extended straightforwardly to the case of Hard Spheres; indeed at some stage is was assumed that the vibrations around the equilibrium positions were harmonic in a first approximation. This approximation is not bad for soft potentials, but it clearly makes no sense for hard spheres. A related but different approach was used in CFP98 , obtaining a reasonable estimate of the Kauzmann density $`\phi _K0.62`$; however, the estimate of the configurational entropy was wrong by two orders of magnitude and the thermodynamic properties of the glass could not be computed within this approach.
The aim of this work is to adapt the replica method of MP99 to the case of the Hard Sphere liquid, and in general of potentials such that the pair distribution function $`g(r)`$ shows discontinuities. This allows us to compute from first principles the configurational entropy of the liquid as well as the thermodynamic properties of the glass and the random close packing density. We find a very good estimate of the configurational entropy that agrees well with recent numerical simulations Sp98 ; Luca05 , a Kauzmann density in the range $`0.58รท0.62`$ (depending on the equation of state we use to describe the liquid state), and a random close packing density in the range $`0.64รท0.67`$. Moreover, we find that the mean coordination number in the amorphous packed states is $`z=6`$ irrespective of the equation of state we use for the liquid, in very good agreement with the result of numerical simulations Be72 ; Ma74 ; SEGHL02 .
The structure of the paper is the following: in section II we outline the replica method of MP99 ; in section III we show how it can be adapted to the case of Hard Spheres; in section IV we resume the main formulae from which we derive our results; in section V we present our main results about the configurational entropy of the liquid and the thermodynamic properties of the glass; in section VI we discuss the behavior of the correlation functions in the glass phase; finally, in section VII we compare our results with previous works.
## II The replica approach to the structural glass transition
The replica method was successfully adapted to the study of the glass transition of simple liquids in a series of recent papers Mo95 ; MP99 ; MP99b ; MP00 ; CMPV99 . The strategy as well as the physics beyond it have been described in detail in MP99 : in this section we will only review the main steps of this approach in order to establish some notations.
### II.1 The molecular liquid
Let us consider here a system at fixed density as in MP99 . The discussion is trivially extended to the case of interest here where the density is the control parameter.
Close to the glass transition the phase space is disconnected in an exponential number of states. The number of states of free energy $`f`$ is called $`๐ฉ(f)=\mathrm{exp}N\mathrm{\Sigma }(f)`$. The complexity $`\mathrm{\Sigma }(f)`$ is a concave function of $`f`$ and vanishes at some value $`f_{min}`$. One can write the partition function $`Z`$ in the following way:
$$\begin{array}{cc}\hfill Z& =e^{\beta NF(T)}\underset{\alpha }{}e^{\beta Nf_\alpha }\hfill \\ & =_{f_{min}}^{f_{max}}๐fe^{N[\mathrm{\Sigma }(f)\beta f]}e^{N[\mathrm{\Sigma }(f^{})\beta f^{}]},\hfill \end{array}$$
(1)
where $`f^{}`$ is such that $`\beta \mathrm{\Phi }(f)=\beta f\mathrm{\Sigma }(f)`$ is minimum. The ideal glass transition is met at the temperature $`T_K`$ such that $`f^{}(T_K)=f_{min}`$ and $`\mathrm{\Sigma }(f^{})=0`$.
The basic idea of the replica approach Mo95 ; MP99 is to consider $`m`$ copies of the original system, constrained to be in the same state by a small attractive coupling. The partition function of the replicated system is then
$$\begin{array}{cc}\hfill Z_m& =e^{\beta N\mathrm{\Phi }(m,T)}\underset{\alpha }{}e^{\beta Nmf_\alpha }\hfill \\ & =_{f_{min}}^{f_{max}}๐fe^{N[\mathrm{\Sigma }(f)\beta mf]}e^{N[\mathrm{\Sigma }(f^{})\beta mf^{}]},\hfill \end{array}$$
(2)
where now $`f^{}(m,T)`$ is such that $`\beta \mathrm{\Phi }(m,f)=\beta mf\mathrm{\Sigma }(f)`$ is minimum. If $`m`$ is allowed to assume real values, the complexity can be estimated from the knowledge of the function $`\beta \mathrm{\Phi }(m,T)=\beta mf^{}(m,T)\mathrm{\Sigma }(f^{}(m,T))`$. Indeed, it is easy to show that
$$\begin{array}{cc}\hfill \beta f^{}(m,T)& =\frac{\beta \mathrm{\Phi }(m,T)}{m},\hfill \\ \hfill \mathrm{\Sigma }(m,T)& =\mathrm{\Sigma }(f^{}(m,T))=m^2\frac{[m^1\beta \mathrm{\Phi }(m,T)]}{m}\hfill \\ & =m\beta f^{}(m,T)\beta \mathrm{\Phi }(m,T).\hfill \end{array}$$
(3)
The function $`\mathrm{\Sigma }(f)`$ can be reconstructed from the parametric plot of $`f^{}(m,T)`$ and $`\mathrm{\Sigma }(m,T)`$.
Moreover, at fixed $`m<1`$, the glass transition is shifted towards lower values of the temperature. Indeed, for any value of the temperature $`T`$ below $`T_K`$ it exists a value $`m^{}(T)<1`$ such that for $`m<m^{}`$ the system is in the liquid phase. The free energy for $`T<T_K`$ and $`m<m^{}(T)`$ can be computed by analytic continuation of the free energy of the high temperature liquid. As the free energy is always continuous and it is independent of $`m<m^{}(T)`$ in the glass phase (being simply the value $`f_{min}(T)`$ such that $`\mathrm{\Sigma }(f_{min})=0`$), one can compute the free energy of the glass below $`T_K`$ simply as $`F_{glass}(T)=\mathrm{\Phi }(m^{}(T),T)/m^{}(T)`$.
The $`m`$ copies are assumed to be in the same state. This means that each atom of a given replica is close to an atom of each of the other $`m1`$ replicas, i.e., the liquid is made of molecules of $`m`$ atoms, each belonging to a different replica of the original system. In other words the atoms of different replicas stay in the same cage. The replica method allow us to define and compute the properties of the cages in a purely equilibrium framework, in spite of the fact that the cages have been defined originally in a dynamic framework. The problem is then to compute the free energy of a molecular liquid where each molecule is made of $`m`$ atoms. The $`m`$ atoms are kept close one to each other by a small inter-replica coupling that is switched off at the end of the calculation, while each atom interacts with all the other atoms of the same replica via the original pair potential. This problem can be tackled by mean of the HNC integral equations Hansen .
### II.2 HNC free energy
The traditional HNC approximation can be naturally extended to the case where particles have internal degrees of freedom and also to the replica approach where we have molecules composed by $`m`$ atoms.
We will denote by $`x=\{\underset{ยฏ}{x}_1,\mathrm{},\underset{ยฏ}{x}_m\}`$, $`\underset{ยฏ}{x}_a\text{}^d`$ the coordinate of a molecule in dimension $`d`$. The single-molecule density is
$$\rho (x)=\underset{i=1}{\overset{N}{}}\underset{a=1}{\overset{m}{}}\delta (\underset{ยฏ}{x}_{ia}\underset{ยฏ}{x}_a),$$
(4)
and the pair correlation is
$$\rho (x)g(x,y)\rho (y)=\underset{i,j}{\overset{1,N}{}}\underset{a=1}{\overset{m}{}}\delta (\underset{ยฏ}{x}_{ia}\underset{ยฏ}{x}_a)\underset{b=1}{\overset{m}{}}\delta (\underset{ยฏ}{x}_{jb}\underset{ยฏ}{y}_b).$$
(5)
We define also $`h(x,y)=g(x,y)1`$. The interaction potential between two molecules is $`v(x,y)=_av(|\underset{ยฏ}{x}_a\underset{ยฏ}{y}_a|)`$.
The HNC free energy is given by Hansen ; MP99
$$\begin{array}{cc}& \beta \mathrm{\Psi }[\rho (x),g(x,y)]=\frac{1}{2}dxdy\rho (x)\rho (y)[g(x,y)\mathrm{log}g(x,y)\hfill \\ & g(x,y)+1+\beta v(x,y)g(x,y)]\hfill \\ & +๐x\rho (x)\left[\mathrm{log}\rho (x)1\right]+\frac{1}{2}\underset{n3}{}\frac{(1)^n}{n}\text{Tr}[h\rho ]^n,\hfill \end{array}$$
(6)
where
$$\begin{array}{cc}\hfill \text{Tr}[h\rho ]^n& =๐x_1\mathrm{}๐x_nh(x_1,x_2)\rho (x_2)h(x_2,x_3)\rho (x_3)\hfill \\ & \mathrm{}h(x_{n1},x_n)\rho (x_n)h(x_n,x_1)\rho (x_1).\hfill \end{array}$$
(7)
For Hard Spheres the potential term vanishes, $`๐x๐y\rho (x)\rho (y)g(x,y)v(x,y)0`$, so the reduced free energy $`\beta \mathrm{\Psi }`$ will not depend on the temperature in all the following equations. Similarly, all the free energy functions that we will consider below do not depend on the temperature once multiplied by $`\beta `$. In principle we could stick to $`\beta =1`$ and slightly simplify the formulae. We have preferred to keep explicitly $`\beta `$, in order to conform to the standard notation for soft spheres (or for hard spheres with an extra potential).
Differentiation w.r.t $`g(x,y)`$ leads to the HNC equation:
$$\mathrm{log}g(x,y)+\beta v(x,y)=h(x,y)c(x,y),$$
(8)
having defined $`c(x,y)`$ from
$$h(x,y)=c(x,y)+๐zc(x,z)\rho (z)h(z,y).$$
(9)
The free energy (per particle) of the system is given by
$$\begin{array}{cc}& \varphi (m,T)=\frac{1}{Nm}\underset{\rho (x),g(x,y)}{\mathrm{min}}\mathrm{\Psi }[\rho (x),g(x,y)],\hfill \\ & \mathrm{\Phi }(m,T)=m\varphi (m,T),\hfill \end{array}$$
(10)
and once the latter is known one can get the free energy of the states and the complexity using Eq.s (3).
### II.3 Single molecule density
The solution of the previous equations for generic $`m`$ is a very complex problem (it is already rather difficult for $`m=2`$). Some kind of ansatz is needed to simplify the computation, that may become terribly complicated.
The single molecule density encodes the information about the inter-replica coupling that keeps all the replicas in the same state. We assume that this arbitrarily small coupling has already been switched off, with the main effect of building molecules of $`m`$ atoms vibrating around the center of mass $`\underset{ยฏ}{X}\text{}^d`$ of the molecule with a certain โcage radiusโ $`A`$. The simplest ansatz for $`\rho (x)`$ is then MP99
$$\rho (x)=\widehat{\rho }๐\underset{ยฏ}{X}\underset{a}{}\rho (\underset{ยฏ}{x}_a\underset{ยฏ}{X}),๐\underset{ยฏ}{u}\rho (\underset{ยฏ}{u})=1,$$
(11)
with
$$\rho (\underset{ยฏ}{u})=\frac{e^{\frac{u^2}{2A}}}{(\sqrt{2\pi A})^d},$$
(12)
and $`\widehat{\rho }=V^1๐x\rho (x)`$ the number density of molecules. With this choice it is easy to show that
$$\begin{array}{cc}\hfill \frac{1}{N}& dx\rho (x)\left[\mathrm{log}\rho (x)1\right]=\mathrm{log}\widehat{\rho }1+\hfill \\ & \frac{d}{2}(1m)\mathrm{log}(2\pi A)\frac{d}{2}\mathrm{log}m+\frac{d}{2}(1m)\hfill \end{array}$$
(13)
### II.4 Pair correlation
As the information about the inter-replica coupling is already encoded in $`\rho (x)`$, we make the ansatz for $`g(x,y)`$:
$$g(x,y)=\underset{a}{}g(|\underset{ยฏ}{x}_a\underset{ยฏ}{y}_a|),$$
(14)
where $`g(r)`$ is rotationally invariant because so is the interaction potential. We also define $`G(r)[g(r)]^m`$. Using the ansatz above, it is easy to rewrite the free energy (6) as follows:
$$\begin{array}{cc}\hfill \beta \mathrm{\Psi }& =\frac{\widehat{\rho }N}{2}d\underset{ยฏ}{r}\{m[F_0(r)]^{m1}F_1(r)[F_0(r)]^m\hfill \\ & +1+m[F_0(r)]^{m1}F_v(r)\}\hfill \\ & +๐x\rho (x)\left[\mathrm{log}\rho (x)1\right]+\frac{1}{2}\underset{n3}{}\frac{(1)^n}{n}\text{Tr}[h\rho ]^n,\hfill \end{array}$$
(15)
where
$$\begin{array}{cc}& F_p(|\underset{ยฏ}{r}|)=๐\underset{ยฏ}{u}๐\underset{ยฏ}{v}\rho (\underset{ยฏ}{u})\rho (\underset{ยฏ}{v})g(|\underset{ยฏ}{r}+\underset{ยฏ}{u}\underset{ยฏ}{v}|)[\mathrm{log}g(|\underset{ยฏ}{r}+\underset{ยฏ}{u}\underset{ยฏ}{v}|)]^p\hfill \\ & F_v(|\underset{ยฏ}{r}|)=๐\underset{ยฏ}{u}๐\underset{ยฏ}{v}\rho (\underset{ยฏ}{u})\rho (\underset{ยฏ}{v})g(|\underset{ยฏ}{r}+\underset{ยฏ}{u}\underset{ยฏ}{v}|)\beta v(|\underset{ยฏ}{r}+\underset{ยฏ}{u}\underset{ยฏ}{v}|)\hfill \end{array}$$
(16)
Note that as $`g(r)`$ and $`v(r)`$ are rotationally invariant, so are $`F_p(r)`$ and $`F_v(r)`$. If $`\rho (\underset{ยฏ}{u})`$ is given by Eq. (12), one gets
$$F(|\underset{ยฏ}{r}|)=๐\underset{ยฏ}{u}\frac{e^{\frac{u^2}{4A}}}{(\sqrt{4\pi A})^d}f(|\underset{ยฏ}{r}+\underset{ยฏ}{u}|)$$
(17)
where $`f(r)\{g(r),g(r)\mathrm{log}g(r),g(r)\beta v(r)\}`$. For Hard Spheres $`F_v0`$.
## III Small cage expansion
The strategy of MP99 was to expand the HNC free energy in a power series of the cage radius $`A`$, assuming that the latter is small close to the glass transition. The expansion is carried out easily if the pair potential $`v(r)`$ and the pair correlation $`g(r)`$ are analytic functions of $`r`$. However this is not the case for Hard Spheres, as $`g(r)`$ vanishes for $`r<D`$ and has a discontinuity in $`r=D`$, so the formulae of MP99 for the power series expansion of $`\mathrm{\Psi }`$ cannot be applied to our system. In this section, we will work out the expansion in the case where the pair correlation $`g(r)`$ has discontinuities.
It is crucial to realize, that independently from any approximation, in the limit $`A0`$, the partition function becomes (neglecting a trivial factor) the partition function of a single atom at an effective temperature given by $`\beta _{eff}=\beta m`$. In the case of hard spheres, where there is no dependence on the temperature, the change in temperature is irrelevant.
In MP99 it was shown that the first term of the expansion is proportional to $`A`$ if $`g(r)`$ is differentiable. As we will see in the following, in the case of hard spheres, the presence of a jump in $`g(r)`$ produces terms $`O(\sqrt{A})`$ in the expansion. In this paper we will focus on these terms neglecting all the contributions of higher order in $`\sqrt{A}`$. This means that we can neglect all the contributions coming from the regions where $`g(r)`$ is differentiable and concentrate only on what happens around $`r=D`$.
We will focus first on the $`g(\mathrm{log}g1)`$ term in Eq. (6). The contribution we want to estimate comes from the discontinuity of $`g(r)`$ in $`r=D`$. Thus to compute this correction the form of $`g(r)`$ away from the singularity is irrelevant and we will use the simplest possible form of $`g(r)`$.
### III.1 Expansion of $`F_0(r)`$
First we will discuss the expansion of $`F_0(r)`$ in $`d=1`$. The simplest possible form of $`g(r)`$ is
$$g(r)=\theta (rD)[1+(y1)e^{\mu (rD)}];$$
(18)
the amplitude of the jump of $`g(r)`$ in $`r=D`$ is given by $`y`$. Remember that in our notation $`\underset{ยฏ}{r}\text{}`$ and $`r=|\underset{ยฏ}{r}|\text{}^+`$. As the functions $`F_0`$ and $`g`$ are even in $`\underset{ยฏ}{r}`$, we can write
$$_{\mathrm{}}^{\mathrm{}}๐\underset{ยฏ}{r}[F_0(\underset{ยฏ}{r})^mg(\underset{ยฏ}{r})^m]=2_0^{\mathrm{}}๐r[F_0(r)^mg(r)^m].$$
(19)
Defining
$$\begin{array}{cc}& \text{erf}(t)\frac{2}{\sqrt{\pi }}_0^t๐xe^{x^2},\hfill \\ & \mathrm{\Theta }(t)=\frac{1}{2}[1+\text{erf}(t)],\hfill \end{array}$$
(20)
these functions play the role of โsmoothedโ sign and $`\theta `$-function respectively; note also that the function $`\mathrm{\Theta }(t)`$ goes to $`0`$ as $`e^{t^2}`$ for $`t\mathrm{}`$. Then
$$\begin{array}{cc}& _{\mathrm{}}^{\mathrm{}}๐u\frac{e^{\frac{u^2}{4A}}}{\sqrt{4\pi A}}\theta (r+uD)=\hfill \\ & \frac{1}{2}\left[1+\text{erf}\left(\frac{rD}{\sqrt{4A}}\right)\right]\mathrm{\Theta }\left(\frac{rD}{\sqrt{4A}}\right),\hfill \end{array}$$
(21)
and
$$\begin{array}{cc}\hfill F_0(r)& =\mathrm{\Theta }\left(\frac{rD}{\sqrt{4A}}\right)+\mathrm{\Theta }\left(\frac{r+D}{\sqrt{4A}}\right)\hfill \\ & +(y1)e^{A\mu ^2}\{e^{\mu (rD)}\mathrm{\Theta }\left(\frac{rD2A\mu }{\sqrt{4A}}\right)\hfill \\ & +e^{\mu (r+D)}\mathrm{\Theta }(\frac{r+D+2A\mu }{\sqrt{4A}})\}.\hfill \end{array}$$
(22)
As $`r0`$ we can neglect the terms proportional to $`\mathrm{\Theta }\left(\frac{r+D}{\sqrt{4A}}\right)`$ in Eq. (22), that give a contribution of order $`\mathrm{exp}(D^2/A)`$ for $`A0`$. Defining the reduced variable $`t=(rD)/\sqrt{4A}`$:
$$\begin{array}{cc}\hfill g(t)& =\theta (t)[1+(y1)e^{\mu 2\sqrt{A}t}],\hfill \\ \hfill F_0(t)& =\mathrm{\Theta }(t)+(y1)e^{\mu 2\sqrt{A}t}e^{A\mu ^2}\mathrm{\Theta }(t+\mu \sqrt{A}),\hfill \end{array}$$
(23)
and Eq. (19) becomes
$$\begin{array}{cc}& _0^{\mathrm{}}๐r[F_0(r)^mg(r)^m]=\hfill \\ & 2\sqrt{A}_{\frac{D}{\sqrt{4A}}}^{\mathrm{}}๐t[F_0(t)^mg(t)^m]2\sqrt{A}Q(A).\hfill \end{array}$$
(24)
If the function $`Q(A)`$ has a finite limit $`Q(0)`$ for $`A0`$ we will have $`Q(A)=Q(0)+o(1)`$ and the leading correction to the free energy is $`O(\sqrt{A}Q(0))`$. The limit for $`A0`$ of $`Q(A)`$ is formally given by
$$Q(0)=y^m_{\mathrm{}}^{\mathrm{}}๐t[\mathrm{\Theta }(t)^m\theta (t)^m]y^mQ_m$$
(25)
where $`y^mY`$ is the jump of $`G(r)g(r)^m`$ in $`r=D`$ and $`Q_m_{\mathrm{}}^{\mathrm{}}๐t[\mathrm{\Theta }(t)^m\theta (t)^m]`$. It is easy to show that $`Q_m`$ is a finite and smooth function of $`m`$ for $`m0`$, that
$$\begin{array}{cc}& Q_m=(1m)Q_0+O[(m1)^2],\hfill \\ & Q_0=_{\mathrm{}}^{\mathrm{}}๐t\mathrm{\Theta }(t)\mathrm{log}\mathrm{\Theta }(t)0.638,\hfill \end{array}$$
(26)
and that $`Q_m`$ diverges as $`Q_m\sqrt{\pi /4m}`$ for $`m0`$. Finally we get, recalling that $`G(r)=[g(r)]^m`$,
$$\frac{1}{2}๐\underset{ยฏ}{r}F_0(r)^m=\frac{1}{2}๐\underset{ยฏ}{r}G(r)+2\sqrt{A}YQ_m.$$
(27)
In dimension $`d>1`$ we have, recalling that $`F_0(r)`$ and $`G(r)`$ are both rotationally invariant,
$$๐\underset{ยฏ}{r}[F_0(r)^mG(r)^m]=\mathrm{\Omega }_d_0^{\mathrm{}}๐rr^{d1}[F_0(r)^mG(r)^m],$$
(28)
where $`\mathrm{\Omega }_d`$ is the solid angle in $`d`$ dimension, $`\mathrm{\Omega }_d=2\pi ^{d/2}/\mathrm{\Gamma }(d/2)`$. The function $`F_0(r)`$ can be written as
$$F_0(r)=๐\underset{ยฏ}{u}\frac{e^{\frac{u^2}{4A}}}{(\sqrt{4\pi A})^d}g(|r\widehat{i}+\underset{ยฏ}{u}|),$$
(29)
where $`\widehat{i}`$ is the unit vector e.g. of the first direction in $`\text{}^d`$. For small $`\sqrt{A}`$, the $`u`$ are small too. The function $`g(|r\widehat{i}+\underset{ยฏ}{u}|)`$ is differentiable along the directions orthogonal to $`\widehat{i}`$. Expanding in series of $`u_\mu `$, $`\mu 1`$, at fixed $`u_1`$, we see that the integration over these variables gives a contribution $`O(A)`$, so we finally get:
$$F_0(r)=_{\mathrm{}}^{\mathrm{}}๐u_1\frac{e^{\frac{u_1^2}{4A}}}{\sqrt{4\pi A}}g(r+u_1)+O(A),$$
(30)
as in the one dimensional case. The function $`F_0(r)^mG(r)^m`$ is large only for $`rD\sqrt{A}`$ so at the lowest order we can replace $`r^{d1}`$ with $`D^{d1}`$ in Eq. (28). We get
$$๐\underset{ยฏ}{r}[F_0(r)^mG(r)^m]=\mathrm{\Omega }_dD^{d1}_0^{\mathrm{}}๐r[F_0(r)^mG(r)^m].$$
(31)
The last integral, with $`F_0(r)`$ given by Eq. (30) is the same as in $`d=1`$, so we obtain
$$\frac{1}{2}๐\underset{ยฏ}{r}F_0(r)^m=\frac{1}{2}๐\underset{ยฏ}{r}G(r)+\sqrt{A}Y\mathrm{\Sigma }_d(D)Q_m,$$
(32)
where $`\mathrm{\Sigma }_d(D)`$ is the surface of a $`d`$-dimensional sphere of radius $`D`$, $`\mathrm{\Sigma }_d(D)=\mathrm{\Omega }_dD^{d1}`$. This result can be formally written as
$$\begin{array}{cc}\hfill F_0(r)^m& G(r)+2\sqrt{A}YQ_m\delta (|r|D)\hfill \\ & G(r)+Q_0(r)\hfill \end{array}$$
(33)
as the correction comes only from the region close to the singularity of $`g(r)`$, $`rD\sqrt{A}`$.
### III.2 $`G\mathrm{log}G`$-term
Let us now estimate the correction coming from the term $`๐rmF_0(r)^{m1}F_1(r)`$. Using the same argument as in the previous subsection, we will restrict to $`d=1`$. Note first that $`F_0(r)`$, for $`|rD|\sqrt{A}`$, has the form
$$F_0(r)=y\mathrm{\Theta }\left(\frac{rD}{\sqrt{4A}}\right)+o(\sqrt{A}),$$
(34)
where $`y`$ is the jump of the function $`g(r)`$ in $`r=D`$. Similarly, $`F_1(r)`$ will have the form
$$F_1(r)=\{\begin{array}{cc}g(r)\mathrm{log}g(r)+O(A),|rD|\sqrt{A},\hfill & \\ y\mathrm{log}y\mathrm{\Theta }\left(\frac{rD}{\sqrt{4A}}\right)+o(\sqrt{A}),|rD|\sqrt{A}.\hfill & \end{array}$$
(35)
The integral
$$_0^{\mathrm{}}๐r[mF_0(r)^{m1}F_1(r)mg(r)^m\mathrm{log}g(r)]$$
(36)
has then two contributions: the first comes from the region $`|rD|\sqrt{A}`$ and is of order $`A`$ as if the function $`g(r)`$ were continuous. The other comes from the region $`|rD|\sqrt{A}`$ and is of order $`\sqrt{A}`$ as in the previous case. To estimate the latter we can use again the reduced variable $`t`$ and approximate $`F_1(t)y\mathrm{log}y\mathrm{\Theta }(t)`$, $`F_0(t)y\mathrm{\Theta }(t)`$. Then we get
$$\begin{array}{cc}\hfill _0^{\mathrm{}}๐r& [mF_0(r)^{m1}F_1(r)mg(r)^m\mathrm{log}g(r)]=\hfill \\ & Y\mathrm{log}Y\mathrm{\hspace{0.17em}2}\sqrt{A}Q_m+o(\sqrt{A}),\hfill \end{array}$$
(37)
in $`d=1`$ and finally, in any dimension $`d`$,
$$\begin{array}{cc}& \frac{1}{2}๐\underset{ยฏ}{r}mF_0(r)^{m1}F_1(r)=\hfill \\ & \frac{1}{2}๐\underset{ยฏ}{r}G(r)\mathrm{log}G(r)+\sqrt{A}Y\mathrm{log}Y\mathrm{\Sigma }_d(D)Q_m.\hfill \end{array}$$
(38)
### III.3 Interaction term
Substituting Eq. (11) in the last term of the HNC free energy one obtains
$$\begin{array}{cc}\hfill \text{Tr}[h\rho ]^n& =\widehat{\rho }^nd\underset{ยฏ}{X}_1\mathrm{}d\underset{ยฏ}{X}_ndu_1\mathrm{}du_n\times \hfill \\ & \times \rho (u_1)\mathrm{}\rho (u_n)h(\underset{ยฏ}{X}_1\underset{ยฏ}{X}_2,u_1u_2)\hfill \\ & \mathrm{}h(\underset{ยฏ}{X}_n\underset{ยฏ}{X}_1,u_nu_1),\hfill \end{array}$$
(39)
where we used the notations $`h(X,u)=_{a=1}^mg(X+u_a)1`$ and $`\rho (u)=_{a=1}^m\rho (\underset{ยฏ}{u}_a)`$ with $`\rho (\underset{ยฏ}{u})`$ given by Eq. (12).
The correction $`O(\sqrt{A})`$ to this integral comes from the regions where $`|X_iX_{i+1}|=D+O(\sqrt{A})`$ for some $`i=1,\mathrm{},n`$. In these regions the functions $`h`$ such that their arguments are not close to the singularity can be expanded in a power series in $`u`$, the correction being $`O(A)`$ MP99 . Thus we can write, defining $`H(r)=G(r)1`$:
$$\begin{array}{cc}\hfill \widehat{\rho }^n& \text{Tr}[h\rho ]^n=๐\underset{ยฏ}{X}_1\mathrm{}๐\underset{ยฏ}{X}_nH(\underset{ยฏ}{X}_1\underset{ยฏ}{X}_2)\mathrm{}H(\underset{ยฏ}{X}_n\underset{ยฏ}{X}_1)+\hfill \\ & nd\underset{ยฏ}{X}_1\mathrm{}d\underset{ยฏ}{X}_ndu_1du_2\rho (u_1)\rho (u_2)\times \hfill \\ & \times [h(\underset{ยฏ}{X}_1\underset{ยฏ}{X}_2,u_1u_2)H(\underset{ยฏ}{X}_1\underset{ยฏ}{X}_2)]\times \hfill \\ & \times H(\underset{ยฏ}{X}_2\underset{ยฏ}{X}_3)\mathrm{}H(\underset{ยฏ}{X}_n\underset{ยฏ}{X}_1)=\hfill \\ & ๐\underset{ยฏ}{X}_1\mathrm{}๐\underset{ยฏ}{X}_nH(\underset{ยฏ}{X}_1\underset{ยฏ}{X}_2)\mathrm{}H(\underset{ยฏ}{X}_n\underset{ยฏ}{X}_1)\hfill \\ & +nd\underset{ยฏ}{X}_1\mathrm{}d\underset{ยฏ}{X}_nQ_0(\underset{ยฏ}{X}_1\underset{ยฏ}{X}_2)\times \hfill \\ & \times H(\underset{ยฏ}{X}_2\underset{ยฏ}{X}_3)\mathrm{}H(\underset{ยฏ}{X}_n\underset{ยฏ}{X}_1),\hfill \end{array}$$
(40)
where in the last step we used Eq. (33):
$$\begin{array}{cc}\hfill ๐u_1๐u_2& \rho (u_1)\rho (u_2)\left[h(r,u_1u_2)H(r)\right]=\hfill \\ & F_0(r)^mG(r)=Q_0(r).\hfill \end{array}$$
(41)
Collecting all the terms with different $`n`$ we get
$$\begin{array}{cc}& \frac{1}{2}\underset{n3}{}\frac{(1)^n}{n}\text{Tr}[h\rho ]^n\frac{1}{2}\underset{n3}{}\frac{(1)^n}{n}\widehat{\rho }^n\text{Tr}H^n+\hfill \\ & +\frac{\widehat{\rho }^3}{2}d\underset{ยฏ}{X}_1d\underset{ยฏ}{X}_2d\underset{ยฏ}{X}_3Q_0(\underset{ยฏ}{X}_1\underset{ยฏ}{X}_2)H(\underset{ยฏ}{X}_2\underset{ยฏ}{X}_3)\times \hfill \\ & \times \underset{n3}{}(1)^n\widehat{\rho }^{n3}d\underset{ยฏ}{X}_4\mathrm{}d\underset{ยฏ}{X}_nH(\underset{ยฏ}{X}_3\underset{ยฏ}{X}_4)\times \hfill \\ & \times \mathrm{}H(\underset{ยฏ}{X}_n\underset{ยฏ}{X}_1)=\hfill \\ & =\frac{1}{2}\underset{n3}{}\frac{(1)^n}{n}\widehat{\rho }^n\text{Tr}H^n\frac{\widehat{\rho }^3}{2}d\underset{ยฏ}{X}_1d\underset{ยฏ}{X}_2d\underset{ยฏ}{X}_3\times \hfill \\ & \times Q_0(\underset{ยฏ}{X}_1\underset{ยฏ}{X}_2)H(\underset{ยฏ}{X}_2\underset{ยฏ}{X}_3)C(\underset{ยฏ}{X}_3\underset{ยฏ}{X}_1).\hfill \end{array}$$
(42)
Substituting the expression of $`Q_0(r)`$ and recalling that from the definition of $`C(\underset{ยฏ}{X})`$ one has $`\widehat{\rho }๐\underset{ยฏ}{Z}H(\underset{ยฏ}{X}\underset{ยฏ}{Z})C(\underset{ยฏ}{Z}\underset{ยฏ}{Y})=H(\underset{ยฏ}{X}\underset{ยฏ}{Y})C(\underset{ยฏ}{X}\underset{ยฏ}{Y})`$, we get
$$\begin{array}{cc}& \frac{1}{2}\underset{n3}{}\frac{(1)^n}{n}\text{Tr}[h\rho ]^n=\frac{1}{2}\underset{n3}{}\frac{(1)^n}{n}\widehat{\rho }^n\text{Tr}H^n\hfill \\ & N\widehat{\rho }Q_m\sqrt{A}y\mathrm{\Sigma }_d(D)[H(D)C(D)].\hfill \end{array}$$
(43)
This result is correct in any dimension $`d`$.
## IV First order free energy
Substituting Eq.s (32), (38) and (43) in Eq. (15) one obtains the following expression for the HNC free energy at first order in $`\sqrt{A}`$:
$$\begin{array}{cc}\hfill \beta & F=\frac{\beta \mathrm{\Psi }}{N}=\beta F_0(A)+\beta F_{eq}[G(r)]+\beta \mathrm{\Delta }F[A,G(r)],\hfill \\ \hfill \beta & F_{eq}=\frac{\widehat{\rho }}{2}d^dr\{G(r)\mathrm{log}G(r)G(r)+1\}\hfill \\ & +\frac{1}{2\widehat{\rho }}\frac{d^dk}{(2\pi )^d}\left[\mathrm{log}[1+\widehat{H}(k)]+\widehat{H}(k)\frac{1}{2}\widehat{H}(k)^2\right]\hfill \\ & +\mathrm{log}\widehat{\rho }1,\hfill \\ \hfill \beta & F_0=\frac{d}{2}(1m)\mathrm{log}(2\pi A)+\frac{d}{2}(1m)\frac{d}{2}\mathrm{log}m,\hfill \\ \hfill \beta & \mathrm{\Delta }F=\widehat{\rho }Q_m\sqrt{A}\mathrm{\Sigma }_d(D)G(D)\times \hfill \\ & \times [\mathrm{log}G(D)1H(D)+C(D))],\hfill \end{array}$$
(44)
where $`Q_m=Q_0(1m)+o((m1)^2)`$, $`Q_00.638`$ and the Fourier transform has been defined as
$$\widehat{H}(k)=\widehat{\rho }๐re^{ikr}H(r).$$
(45)
At the first order in $`\sqrt{A}`$ we only need to know the function $`G(r)`$ determined by the optimization of the free energy at the zeroth order in $`\sqrt{A}`$, i.e. the usual free energy $`F_{eq}[G(r)]`$: it satisfies the HNC equation $`\mathrm{log}G(r)=H(r)C(r)`$. Substituting this relation in $`\beta \mathrm{\Delta }F`$ one simply obtains $`\beta \mathrm{\Delta }F=\widehat{\rho }Q_m\sqrt{A}\mathrm{\Sigma }_d(D)G(D)`$.
The derivative w.r.t. $`A`$ leads to the following expression for the cage radius:
$$\sqrt{A^{}}=\frac{1m}{Q_m}\frac{d}{\widehat{\rho }\mathrm{\Sigma }_d(D)G(D)}$$
(46)
which in $`d=3`$ becomes (let us define again $`Y=G(D)`$):
$$\frac{\sqrt{A^{}}}{D}=\frac{1m}{Q_m}\frac{1}{8\phi Y(\phi )}$$
(47)
where $`\phi =\frac{\pi D^3\widehat{\rho }}{6}`$ is the packing fraction. Substituting this result in $`\beta \mathrm{\Delta }F`$ one has $`\beta \mathrm{\Delta }F(A^{})=d(m1)`$.
Finally, the expression for the replicated free energy in $`d=3`$ is
$$\begin{array}{cc}\hfill \beta \mathrm{\Phi }(m,\phi )& =\beta F_{eq}(\phi )+\frac{3}{2}(1m)\mathrm{log}[2\pi A^{}(m)]\hfill \\ & +\frac{3}{2}(m1)\frac{3}{2}\mathrm{log}m\hfill \end{array}$$
(48)
Note that for Hard Spheres one has $`\beta F_{eq}(\phi )=S(\phi )`$, $`S`$ being the total entropy of the liquid. We get then
$$\begin{array}{cc}\hfill \beta f^{}(& m,\phi )=\frac{\beta \mathrm{\Phi }}{m}=\frac{3}{2}\mathrm{log}[2\pi A^{}(m)]\hfill \\ & +\frac{3}{2}(1m)\frac{d\mathrm{log}A^{}(m)}{dm}+\frac{3}{2}\frac{m1}{m},\hfill \\ \hfill \mathrm{\Sigma }(m& ,\phi )=m\beta f^{}\beta \mathrm{\Phi }=S(\phi )\frac{3}{2}\mathrm{log}[2\pi A^{}(m)]\hfill \\ & +\frac{3m}{2}(1m)\frac{d\mathrm{log}A^{}(m)}{dm}+\frac{3}{2}\mathrm{log}m\hfill \end{array}$$
(49)
For small enough density the system is in the liquid phase and $`m`$ is equal to 1 at the saddle point. For $`m=1`$ we have:
$$\begin{array}{cc}& \frac{\sqrt{A^{}(1)}}{D}=\frac{1}{8Q_0\phi Y(\phi )}\hfill \\ & S_{vib}(\phi )\beta f^{}(1,\phi )=\frac{3}{2}\mathrm{log}[2\pi A^{}(1)]\hfill \\ & \mathrm{\Sigma }(\phi )=S(\phi )S_{vib}(\phi )\hfill \end{array}$$
(50)
This allows for a computation of $`\mathrm{\Sigma }(\phi )`$ once $`S(\phi )`$ and $`Y(\phi )`$ are known. Note that $`1+4\phi Y(\phi )=\beta P/\rho =\phi \frac{S}{\phi }`$, so a model for $`S(\phi )`$ (or $`Y(\phi )`$) is enough to determine all the quantities of interest.
## V Results from the HNC free energy
We computed numerically $`S(\phi )`$ and $`Y(\phi )`$ solving the classical HNC equation for the Hard Sphere liquid up to $`\phi =0.65`$. This allows to compute $`\beta \mathrm{\Phi }(\phi ,m)`$ and gives access to all the thermodynamic quantities using Eq.s (49) and (50). In this section we discuss the results of this computation. We will set the sphere diameter $`D=1`$ in the following.
### V.1 Equilibrium complexity
The equilibrium complexity $`\mathrm{\Sigma }(\phi )`$ is given by Eq. (50). It is reported in Fig. 1. We get a complexity $`\mathrm{\Sigma }1`$ as found in previous calculations in Lennard-Jones systems MP99 ; CMPV99 ; MP99b ; MP00 , as well as in the numerical simulations CMPV99 ; SKT99 . The complexity vanishes at $`\phi _K=0.582`$, that is the ideal glass transition density โor Kauzmann densityโ predicted by the HNC equations.
### V.2 Phase diagram in the $`(\phi ,m)`$ plane
We now compute the thermodynamic properties of the glassy phase for $`\phi >\phi _K`$. As discussed above, it exists a value of $`m`$, $`m^{}(\phi )`$, such that for $`m<m^{}(\phi )`$ the system is in the liquid phase. It is the solution of $`\mathrm{\Sigma }(m,\phi )=0`$, where $`\mathrm{\Sigma }(m,\phi )`$ is given by Eq. (49). In Fig. 2 we report $`m^{}`$ as a function of $`\phi `$. Clearly, $`m^{}=1`$ at $`\phi =\phi _K`$ and $`m^{}<1`$ for $`\phi >\phi _K`$. $`m^{}`$ vanishes linearly at $`\phi _c=0.640`$. As we will see in the following, above this value of $`\phi `$ the glassy state does not exist anymore.
### V.3 Thermodynamic properties of the glass
The knowledge of the function $`m^{}(\phi )`$ allows to compute the entropy of the glass. Indeed, the free energy does not depend on $`m`$ in the whole glassy phase, and it is continuous along the line $`m=m^{}(\phi )`$, so we can compute the entropy of the glass simply as
$$S_{glass}(\phi )=\beta F_{glass}(\phi )=\frac{\beta \mathrm{\Phi }(m^{}(\phi ),\phi )}{m^{}(\phi )}$$
(51)
This relation is true for $`m^{}<1`$. Below $`\phi _K`$ one has $`m^{}>1`$ and the liquid phase is the stable one. Eq. (51) for $`m^{}>1`$ gives the entropy of the lowest states in the free energy landscape (see below) and can be regarded as the analytic continuation of the glass entropy below $`\phi _K`$. The reader should notice that the glass phase for $`m^{}>1`$ does not have a simple physical meaning and the interesting part of the curves for the glass is in the region $`\phi >\phi _K`$.
In Fig. 3 we report the entropies of the liquid and the glass as functions of the packing fraction. The glass phase becomes stable above $`\phi _K=0.582`$; note that the entropy of the glass is smaller than the entropy of the liquid, i.e. its free energy is bigger than the free energy of the liquid. The same happens also in Lennard-Jones systems and in mean-field spin glass systems. However the physical relevant parts of the curves are the liquid one for $`\phi <\phi _K`$ and the glassy one for $`\phi >\phi _K`$.
The reduced pressure,
$$\frac{\beta P}{\rho }=\phi \frac{S}{\phi },$$
(52)
is reported in Fig. 4.
It is continuous at $`\phi _K`$ and the glass transition is a second order transition from the thermodynamical point of view. Note that the pressure in the glass phase is well described by a power law and it has a simple pole at $`\phi _c`$:
$$\frac{\beta P_{glass}}{\rho }\frac{1}{\phi _c\phi },$$
(53)
as one can see from the inset of Fig. 4 where the inverse reduced pressure is plotted as a function of $`\phi `$.
For $`\phi \phi _c`$ the pressure of the glass diverges and its compressibility $`\chi =\frac{1}{\phi }\frac{\phi }{P}`$ vanishes and consequently $`\phi _c`$ is the maximum density allowed for a disordered state, i.e. it can be identified as the random close packing density. The value $`\phi _c=0.640`$ is in very good agreement with the values reported in the literature. Note that the compressibility jumps downward on increasing $`\phi `$ across $`\phi _K`$, i.e. the compressibility of the glass is smaller than the compressibility of the liquid.
### V.4 Cage radius
The cage radius is given as a function of $`m`$ in Eq. (47). In Fig. 5 we report the cage radius in the liquid phase, $`\sqrt{A^{}(1)}`$, see Eq. (50), and the cage radius in the glass phase, defined as $`\sqrt{A^{}(m^{})}`$. As $`Q_m\sqrt{\pi /4m}`$ for $`m0`$, the cage radius vanishes as $`\sqrt{m^{}}`$ for $`m^{}0`$, i.e. it is proportional to $`\sqrt{\phi _c\phi }`$. The vanishing of the cage radius for $`\phi \phi _c`$ means that at $`\phi _c`$ each sphere is in contact with its neighbors, that is consistent with our interpretation of $`\phi _c`$ as the random close packing density.
### V.5 Complexity of the metastable states
From the parametric plot of $`\beta f^{}(m,\phi )`$ and $`\mathrm{\Sigma }(m,\phi )`$ given in Eq. (49) by varying $`m`$, one can reconstruct the function $`\mathrm{\Sigma }(\beta f)`$ for each value of the packing fraction. This function is reported in Fig. 6 for some values of $`\phi `$ below and above $`\phi _K`$. The function $`\mathrm{\Sigma }(\beta f)`$ vanishes at a certain value $`\beta f_{min}`$, that is given by Eq. (51). The saddle-point equation that determines the free energy of the equilibrium states is, from Eq. (1),
$$\frac{d\mathrm{\Sigma }(\beta f)}{d\beta f}=1.$$
(54)
From Fig. 6 we see that this equation has a solution $`f^{}>f_{min}`$ for $`\phi <\phi _K=0.582`$. Above $`\phi _K`$ Eq. (54) does not have a solution so the saddle point is simply $`f^{}=f_{min}`$ and the systems goes in the glass state. In this sense, the free energy $`f_{min}`$ of the lowest states below $`\phi _K`$ can be regarded as the analytic continuation of the free energy of the glass, see Fig. 3. The curves $`\mathrm{\Sigma }(\beta f)`$ in Fig. 6 have been truncated arbitrarily at high $`\beta f`$. We have not done consistency checks to investigate where the higher free energy states become unstable (i.e. , to compute $`f_{max}`$).
## VI Correlation functions
We will now turn to the study of the pair distribution function $`\stackrel{~}{g}(r)`$ in the glass state. In principle a full computation would require the evaluation of the corrections proportional to $`\sqrt{A}`$ in the correlation functions of a molecule. However we neglect these terms, that we believe are small, and we consider again our simple ansatz (11), (14) for the correlation function of the molecules, in which the information on the shape of the molecule is only encoded in the function $`\rho (x)`$; these corrections should be physically more relevant and interesting.
As we will see in the following, the correlation function of the spheres in the glass is very similar to the one in the liquid but develops an additional strong peak (that becomes a $`\delta `$-function at $`\phi _c`$) around $`r=D`$. The integral of the latter peak is related to the average coordination number of the random close packings.
### VI.1 Expression of $`\stackrel{~}{g}(r)`$ in the glass phase
We assumed the following form for the pair distribution function of the molecular liquid, see Eq.s (11) and (14):
$$\begin{array}{cc}& \rho _2(x,y)=\rho (x)g(x,y)\rho (y)=\hfill \\ & \widehat{\rho }^2๐X๐Y\underset{a=1}{\overset{m}{}}\rho (x_aX)g(|x_ay_a|)\rho (y_aY).\hfill \end{array}$$
(55)
The pair correlation $`\stackrel{~}{g}(r)`$ of a single replica is obtained integrating over the coordinates of all the replicas but one:
$$\stackrel{~}{g}(|x_1y_1|)=\widehat{\rho }^2๐\underset{ยฏ}{x}_2\mathrm{}๐\underset{ยฏ}{x}_m๐\underset{ยฏ}{y}_2\mathrm{}๐\underset{ยฏ}{y}_m\rho _2(x,y).$$
(56)
Using Eq. (55) we get, with some simple changes of variable:
$$\stackrel{~}{g}(r)=g(r)๐\underset{ยฏ}{u}๐\underset{ยฏ}{v}\rho (\underset{ยฏ}{u})\rho (\underset{ยฏ}{v})F_0(|\underset{ยฏ}{r}+\underset{ยฏ}{u}\underset{ยฏ}{v}|)^{m1},$$
(57)
where $`F_0(r)`$ is defined in Eq. (16). The HNC free energy is optimized by $`g(r)=G(r)^{1/m}`$, where $`G(r)`$ is the HNC pair correlation. Thus we get the following expression for the pair correlation of a single replica:
$$\begin{array}{cc}& \stackrel{~}{g}(r)=G(r)^{\frac{1}{m}}๐\underset{ยฏ}{u}\frac{e^{\frac{u^2}{4A}}}{(\sqrt{4\pi A})^d}F_0(|\underset{ยฏ}{r}+\underset{ยฏ}{u}|)^{m1},\hfill \\ & F_0(r)=๐\underset{ยฏ}{u}\frac{e^{\frac{u^2}{4A}}}{(\sqrt{4\pi A})^d}G(|\underset{ยฏ}{r}+\underset{ยฏ}{u}|)^{\frac{1}{m}}.\hfill \end{array}$$
(58)
For $`m=1`$, i.e. in the liquid phase, this function is trivially equal to $`G(r)`$. This is not the case in the glass phase where $`m<1`$.
### VI.2 Small cage expansion of the correlation function
We will now expand Eq. (58) for small $`A`$. Note first that, if $`rD`$, the function $`g(r+u)`$ can be expanded in powers of $`u`$, and the first correction to $`\stackrel{~}{g}(r)`$ is of order $`A`$. Then, as before, we will concentrate on what happens around $`r=D`$. As already discussed in section III, around $`r=D`$ we have, as in Eq. (34), $`G(r)Y\theta (rD)`$ and
$$F_0(r)Y^{\frac{1}{m}}\mathrm{\Theta }\left(\frac{rD}{\sqrt{4A}}\right),$$
(59)
and Eq. (58) becomes
$$\stackrel{~}{g}(r)=Y\theta (rD)๐\underset{ยฏ}{u}\frac{e^{\frac{u^2}{4A}}}{(\sqrt{4\pi A})^d}\mathrm{\Theta }\left(\frac{|\underset{ยฏ}{r}+\underset{ยฏ}{u}|D}{\sqrt{4A}}\right)^{m1}.$$
(60)
Applying the same argument we used in section III when studying the function $`F_0(r)`$ in dimension $`d>1`$, we can show that the integration over the coordinates $`u_\mu `$, $`\mu 1`$, gives a contribution $`O(A)`$. Then we can rewrite, in any dimension $`d`$:
$$\begin{array}{cc}& \stackrel{~}{g}(r)Y\theta (rD)_{\mathrm{}}^{\mathrm{}}๐u\frac{e^{\frac{u^2}{4A}}}{\sqrt{4\pi A}}\mathrm{\Theta }\left(\frac{r+uD}{\sqrt{4A}}\right)^{m1}\hfill \\ & =G(r)\left\{1+_{\mathrm{}}^{\mathrm{}}\frac{dt}{\sqrt{\pi }}e^{\left(\frac{rD}{\sqrt{4A}}t\right)^2}\left[\mathrm{\Theta }(t)^{m1}1\right]\right\},\hfill \end{array}$$
(61)
defining the reduced variable $`t=\frac{r+uD}{\sqrt{4A}}`$. The second term in the latter expression is a contribution localized around $`r=D`$.
### VI.3 Number of contacts
To compute the average number of contacts, let us recall that the average number of particles in a shell $`[r,r+dr]`$, if there is a particle in the origin, is given by
$$dn(r)=\mathrm{\Omega }_dr^{d1}\widehat{\rho }\stackrel{~}{g}(r)dr.$$
(62)
Thus the number of contacts can be obtained from the correlation function $`\stackrel{~}{g}(r)`$. While the full computation of the correlation function is rather involved, here we limit ourselves to consider the second term in Eq. (61), which is proportional to a Gaussian with variance $`O(\sqrt{A})`$ that becomes a $`\delta (|r|D)`$-function in the limit $`A0`$.
The value of the number of spheres in contact with the sphere in the origin is given by
$$z=\mathrm{\Omega }_d\widehat{\rho }_D^{D+O(\sqrt{A})}๐rr^{d1}\stackrel{~}{g}(r).$$
(63)
The first term in Eq. (61) gives a contribution $`O(\sqrt{A})`$ that can be neglected. If we use $`rD`$ and $`G(r)Y`$ at the leading order in $`\sqrt{A}`$ we obtain, defining the variable $`ฯต=\frac{rD}{\sqrt{4A}}`$,
$$\begin{array}{cc}& z=\mathrm{\Omega }_dD^{d1}\widehat{\rho }Y\times \hfill \\ & \times \sqrt{4A}_0^{\mathrm{}}๐ฯต_{\mathrm{}}^{\mathrm{}}\frac{dt}{\sqrt{\pi }}e^{(ฯตt)^2}\left[\mathrm{\Theta }(t)^{m1}1\right].\hfill \end{array}$$
(64)
Recalling that
$$\frac{1}{\sqrt{\pi }}_0^{\mathrm{}}๐ฯตe^{(ฯตt)^2}=\mathrm{\Theta }(t),$$
(65)
we get, observing that $`_{\mathrm{}}^{\mathrm{}}๐t\left[\mathrm{\Theta }(t)\theta (t)\right]=0`$, and using Eq. (46),
$$\begin{array}{cc}\hfill z& =\mathrm{\Sigma }_d(D)\widehat{\rho }Y\sqrt{4A}_{\mathrm{}}^{\mathrm{}}๐t\mathrm{\Theta }(t)\left[\mathrm{\Theta }(t)^{m1}1\right]\hfill \\ & =\mathrm{\Sigma }_d(D)\widehat{\rho }Y\sqrt{4A}Q_m=2d(1m).\hfill \end{array}$$
(66)
This is the expression of the average number of contacts at the leading order in $`\sqrt{A}`$, to be computed at $`m=m^{}`$ in the glass phase. At $`\phi =\phi _c`$, where $`m^{}=0`$, each sphere has on average $`2d`$ contacts. This is exactly what is found in numerical simulations; the condition $`z2d`$ is required for the mechanical stability of the packings as can be understood by mean of a very simple argument Al98 .
Note that this result is independent on the particular expression we chose for $`S(\phi )`$, $`Y(\phi )`$ and $`G(r)`$, i.e. it might hold beyond the choice of HNC equations for the molecular liquid provided that the expression (46) for the cage radius is correct.
## VII Discussion
We will now compare our results with related ones that appeared in the literature. The main obstacle for a quantitative comparison is that the HNC equations are known to yield a not very good description of the Hard Sphere liquid at high density Hansen ; typically one would obtain the right curves if one shifts the value of $`\phi `$ of a quantity of order 0.03. Therefore, we should limit ourselves to a qualitative comparison of the results coming from the HNC equations with the results of numerical simulations. However, note that, although the expressions (47), (48) for the replicated free energy have been derived starting from the expression (6) for the HNC free energy, the final result depends only on the equilibrium entropy of the liquid $`S(\phi )`$. It is interesting then, for the purpose of comparing our results with experiments and numerical simulations, to consider a more accurate model for $`S(\phi )`$ in the liquid phase. We repeated the calculations of section V substituting the CarnahanโStarling (CS) entropy Hansen
$$\begin{array}{cc}& S_{CS}(\phi )=\mathrm{log}\left(\frac{6\phi }{\pi e}\right)\frac{4\phi 3\phi ^2}{(1\phi )^2},\hfill \\ & Y_{CS}(\phi )=\frac{1\frac{1}{2}\phi }{(1\phi )^3}.\hfill \end{array}$$
(67)
instead of the HNC entropy in Eq.s (48), (47). All the results of section V are qualitatively reproduced using the CS entropy, but the latter gives results in better agreement with the numerical data. However, this procedure is not completely consistent from a theoretical point of view: one should always keep in mind that our aim here is not to present a quantitative theory, but only to show that the replica approach yields a reasonable qualitative scenario for the glass transition in Hard Sphere systems.
### VII.1 Complexity of the liquid and Kauzmann density
In Fig. 7 we report the equilibrium complexity $`\mathrm{\Sigma }(\phi )`$ obtained substituting the HNC and the CS expression for $`S(\phi )`$ and $`Y(\phi )`$ in Eq. (50). The results are compared with recent numerical results of Angelani et al. Luca05 obtained on a $`50:50`$ binary mixture of spheres (to avoid crystallization) with diameter ratio equal to $`1.2`$: the vibrational entropy was estimated using the procedure described in CMPV99 ; AFST04 and the complexity was computed as $`S(\phi )S_{vib}(\phi )`$. A quantitative comparison is difficult here because in the case of a mixture there can be corrections related to the mixing entropy, $`S_{mix}\mathrm{log}2`$. Nevertheless the data are in good agreement with our results. A detailed comparison would require the extension of our computation to binary mixtures following CMPV99 .
Another numerical estimate of $`\mathrm{\Sigma }(\phi )`$ was previously reported by Speedy Sp98 , who rationalized his numerical data assuming a Gaussian distribution of states and a particular form for the vibrational entropy inside a state. The free parameters were then fitted from the liquid equation of state. The curve obtained by Speedy also agrees with our results.
Both the HNC and the CS estimates of the Kauzmann density ($`\phi _K=0.582`$ and $`\phi _K=0.617`$ respectively) fall, as it should be, between the ModeโCoupling dynamical transition that is $`\phi _{MCT}0.56`$ GS91 ; vMU93 , and the Random Close Packing density that is estimated in the range $`\phi =0.64รท0.67`$, see e.g. Be83 .
A computation of $`\mathrm{\Sigma }(\phi )`$ based on very similar ideas was presented in CFP98 , where a very similar estimate of $`\phi _K0.62`$ was obtained. However in CFP98 the complexity was found to be $`\mathrm{\Sigma }0.01`$, i.e. two orders of magnitude smaller than the one obtained from the numerical simulations. This negative result is probably due to some technical problem in the assumptions of CFP98 .
### VII.2 Equation of state of the glass
In Fig. 8 we report as black dots the numerical data for the pressure of the Hard Sphere liquid at high $`\phi `$ obtained by Rintoul and Torquato RT96 . The data were obtained extrapolating at long times the relaxation of the pressure as a function of time after an increase of density starting from an equilibrated configuration at lower density. We also report the curves of the pressure as a function of the density obtained from the HNC and CS equations, both in the liquid and in the glass state.
The agreement of the HNC curve with the data is not very good even in the liquid phase, due to the modest accuracy of the HNC equation of state. However, the qualitative behavior of our curve is in good agreement with the numerical data, and in particular the quasiโlinear behavior of the inverse reduced pressure in the glass phase found in RT96 ; Sp98 , $`\frac{\rho }{\beta P}\phi _c\phi `$, is reproduced by the HNC curve. The HNC pressure of the glass diverges at $`\phi _c=0.640`$ as discussed in section V; the latter is the HNC estimate of the random close packing density.
The CS curve describes well the pressure in the liquid phase Hansen . Comparing the curve with the data of Rintoul and Torquato, we see that the glass transition happens in the numerical simulation at a density $`\phi _g0.56`$ smaller than the one predicted by the CS curve, $`\phi _K=0.617`$ nota1 , and very close to the ModeโCoupling transition density, $`\phi _{MCT}0.56`$. This is not surprising, since the relaxation time grows fast on approaching the ideal glass transition; at some point it becomes larger than the experimental time scale and the liquid falls out of equilibrium becoming a real glass. It is likely that the data of Ref. RT96 describe the pressure of a real nonequilibrium glass, while our computation gives the pressure of the ideal equilibrium glass, that cannot be reached experimentally in finite time.
### VII.3 Random close packing
Both the HNC and CS equations predict the existence of a random close packing density $`\phi _c`$ where the pressure and the value of the radial distribution function $`\stackrel{~}{g}(r)`$ in $`r=D`$ diverge. The HNC estimate is $`\phi _c=0.640`$, in the range of the values ($`\phi _c=0.64รท0.67`$) reported in the literature. The CS estimate is $`\phi _c=0.683`$ and it is also a value consistent with numerical simulations.
The reader should notice that the theoretical value for $`\phi _c`$ is related to the ideal random close packing; however the states corresponding to this value of $`\phi _c`$ can be reached by local algorithms, like most of the algorithms that were used in the literature, in a time that should diverge exponentially with the volume. Some caution should be taken in using the data obtained by numerical simulations. The question of which is the value of the density that can be obtained in large, but finite amount of time per particle is very interesting and more relevant from a practical point of view: however we plan to study it at a later time.
Note that the computation of the mean coordination number $`z`$ of section VI, that gives $`z=6`$ at $`\phi =\phi _c`$ in $`d=3`$, is independent of the particular form we choose for $`S(\phi )`$, and thus is valid for both the HNC and CS equations of state. The value $`z=6`$ has been reported in many studies Be72 ; Ma74 ; Po79 ; Al98 ; SEGHL02 .
## VIII Conclusions
We successfully applied the replica method of Mo95 ; MP99 to the study of the ideal glass transition of Hard Spheres, and in general of potentials such that the pair distribution function $`g(r)`$ shows discontinuities, starting from the replicated HNC free energy and expanding it at first order in the cage radius $`\sqrt{A}`$.
This result allowed us to compute from first principles the configurational entropy of the liquid as well as the thermodynamic properties of the glass up to the random close packing density. Our computation is based on the HNC equation of state, that is known to yield a poor quantitative description of the liquid state at high density. Nevertheless, we found that the qualitative scenario for the ideal glass transition that emerges from the replicated HNC free energy is very reasonable. In particular, we found a complexity $`\mathrm{\Sigma }1`$, a Kauzmann density $`\phi _K=0.582`$, and a random close packing density $`\phi _c=0.64`$. All these results compare well with numerical simulations.
Using, on a phenomenological ground, the CarnahanโStarling equation of state instead of the HNC equation of state as input for our calculations, we could also compare our results with the highโdensity pressure data of Rintoul and Torquato showing that they are indeed compatible with the observation of a real glass transition.
Moreover, we found that the mean coordination number in the amorphous packed states is $`z=2d`$ irrespective of the equation of state we use for the liquid, in very good agreement with the result of numerical simulations and with theoretical arguments Be72 ; Ma74 ; SEGHL02 ; Al98 .
It is worth to note that our results do not prove the existence of a glass transition for the Hard Sphere liquid, as they derive from a particular approximation for the molecular liquid free energy (the HNC approximation), and, in general, other approximation such as the PercusโYevick are possible Hansen .
###### Acknowledgements.
We are grateful to L. Angelani, G. Foffi and F. Sciortino for providing their data prior to publication and for their comments on this work. F.Z. wish also to thank E. Zaccarelli for the code for solving the HNC equations and for many interesting discussions.
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# CHIRAL SOLITON MODEL PREDICTIONS FOR PENTAQUARKS
## 1 Do we see $`\mathrm{\Theta }^+`$ at all?
After almost two year excitement that the exotic antidecuplet has been discovered $`^\mathrm{?}`$ the results from high statistics G11 experiment at CLAS were presented in April at the APS meeting with negative result for the photoproduction of $`\Theta ^+`$ on proton $`^\mathrm{?}`$. The sighting of the heaviest members of $`\overline{10}`$ that were seen only by NA49 experiment at CERN $`^\mathrm{?}`$ is even more problematic. Nevertheless the positive evidence of 11 experiments that reported the existence of $`\mathrm{\Theta }^+`$ cannot be simply ignored. The reasons why some experiments see $`\mathrm{\Theta }^+`$ while the others do not maybe either of experimental nature or a peculiar production mechanism or both. Therefore the present confusion concerning exotics calls for a new high precision $`KN`$ experiment in the interesting energy range.
## 2 Chiral models
Light antidecuplet was predicted within the chiral soliton models ($`\chi `$SM) $`^\mathrm{?}`$<sup>-</sup>$`^\mathrm{?}`$. Early estimate $`\Delta M_{\overline{10}8}600`$ MeV was obtained already in 1984 in a specific modification of the Skyrme model $`^\mathrm{?}`$. The estimates of *both* $`\Theta ^+`$ and $`\Xi _{\overline{10}}`$ masses obtained in the Skyrme model in 1987 are in a surprising agreement with present experimental findings $`^\mathrm{?}`$.
In this Section we will demonstrate that chiral models are deeply rooted in QCD and take into account quark degrees of freedom maybe even in a more complete way than the quark models themselves. The low energy effective theory of QCD could be in principle obtained by integrating out gluons. The resulting quark lagrangian would preserve chiral symmetry, whose spontaneous breakdown would produce nonzero constituent quark mass $`M`$ and the massless pseudoscalar Goldstone bosons, being at the same time $`\overline{\psi }\psi `$ pairs, would be present. A convenient model of such a lagrangian is provided by a semibosonized NambuโJona-Lasinio model:
$$=\overline{\psi }(i\partial ฬธMU^{\gamma _5}[\phi ])\psi $$
(1)
which looks like a Dirac Lagrangian density for a massive fermion $`\psi `$ if not for matrix $`U`$. In fact $`\psi `$ is a $`3`$-vector in flavor space and also in color. Matrix
$$U^{\gamma _5}=\mathrm{exp}\{\frac{i}{F_\phi }\stackrel{}{\lambda }\stackrel{}{\phi }\gamma _5\}$$
(2)
parameterized by a set of eight pseudoscalar fields $`\stackrel{}{\phi }`$ guaranties chiral symmetry of $``$, given by a global multiplication of the fermion field by a phase factor
$$\psi e^{i\stackrel{}{\lambda }\stackrel{}{\alpha }\gamma _5}\psi $$
(3)
provided we also transform meson fields
$$U^{\gamma _5}[\phi ]e^{i\stackrel{}{\lambda }\stackrel{}{\alpha }\gamma _5}U^{\gamma _5}[\phi ]e^{i\stackrel{}{\lambda }\stackrel{}{\alpha }\gamma _5}.$$
(4)
Note that the color indices produce simply an overall factor $`N_c`$ in front of (1).
Since the vacuum state corresponds to $`U^{\gamma _5}=1`$, spontaneous chiral symmetry breaking indeed takes place. Moreover the massless Goldstone bosons appear when we integrate out the quark fields. Then the resulting effective action contains only meson fields and can be organized in terms of a derivative expansion
$$S_{\mathrm{eff}}[\phi ]=\frac{F_\phi ^2}{4}\mathrm{Tr}\left(_\mu U^\mu U^{}\right)+\frac{1}{32e^2}\mathrm{Tr}\left([_\mu UU^{},_\nu UU^{}]^2\right)+\Gamma _{\mathrm{WZ}}+\mathrm{}$$
(5)
where constants $`F_\phi `$ and $`e`$ can be calculated from (1) with an appropriate cut-off. $`\Gamma _{\mathrm{WZ}}`$ is the Witten Wess-Zumino term which takes into account axial anomaly. Perhaps the most important part are the ellipses which encode an infinite set of terms that are effectively summed up by the fermionic model of Eq.(1). The truncated series of Eq.(5) is the basis of the Skyrme model. Hence the Skyrme model is (a somewhat arbitrary, because it does not include another possible 4 derivative term) approximation to (1).
At this point both models, chiral quark model of Eq.(1) and Skyrme model of Eq.(5) (without the โdotsโ), look like mesonic theories describing only meson-meson scattering $`^\mathrm{?}`$. Baryons are introduced in two steps, following large $`N_c`$ strategy described by Witten $`^\mathrm{?}`$. First, one constructs a soliton solution, *i.e.* solution to the classical equations of motion that corresponds to matrix $`U_0`$ which cannot be expanded in a power series around unity. Second, since the classical soliton has no quantum numbers (except baryon number), one has to quantize the system. Perhaps this quantization procedure, which reduces both models to the nonrelativistic quantum system analogous to the symmetric top $`^\mathrm{?}`$ with two moments of inertia $`I_{1,2}`$, makes chiral-soliton models look odd and counterintuitive.
In chiral quark soliton models stabilization of the soliton occurs due to the valence quark level which also provides the baryon number. In the Skyrme model where no quarks are present the soliton is stable due to the specific choice of the 4-derivative term in (5) and the baryon number is given as a charge of the conserved topological current. The quantization on the other hand proceeds in both models almost identically $`^\mathrm{?}`$, the only difference being that some model parameters dominated by the valence level in the quark soliton model are exactly zero in the Skyrme model.
## 3 Exotics in chiral models
Chiral soliton models predict that positive parity baryons fall into SU(3) representations that contain hypercharge $`Y=N_c/3`$ which is 1 in the real world. Therefore the lowest lying multiples are octet and decuplet, exactly as in the quark models. Moreover chiral models predict a tower of exotic rotational states starting with $`\overline{10}_{1/2}`$, $`27_{3/2,1/2}`$, $`\overline{35}_{5/2,3/2}`$ (subscripts refer to spin) etc. The splittings between the centers of the lowest-lying octet, decuplet and antidecuplet baryons are given in the $`\chi `$SM by
$$\mathrm{\Delta }M_{108}=\mathrm{\hspace{0.33em}3}/(2I_1),\mathrm{\Delta }M_{\overline{10}8}=N_c/(2I_2)=\mathrm{\hspace{0.33em}3}/(2I_2)$$
(6)
where $`I_{1,2}`$ are two soliton moments of inertia that depend on details of the chiral Lagrangian. Since $`I_1`$, $`I_2๐ช(N_c)`$, this means that $`\mathrm{\Delta }M_{\overline{10}8}๐ช(N_c^0)`$, whereas $`\mathrm{\Delta }M_{108}`$ is $`๐ช(1/N_c)`$. This has triggered some arguments $`^{\mathrm{?},\mathrm{?}}`$ and counter-arguments $`^\mathrm{?}`$, regarding the applicability of collective coordinate quantization to the $`\overline{10}`$.
We have already mentioned early estimates of the antidecuplet mass that have been recently reviewed in $`^\mathrm{?}`$. The bottom line is that antidecuplet is much lighter than in the quark models. Therefore $`\chi `$SMโs predict light exotic baryons belonging to antideucuplet of positive parity.
Perhaps the most striking prediction of $`\chi `$SM is the small width of the antidecuplet states. The decay width is calculated by means of the formula for the decay width for $`BB^{}+\phi `$:
$$\Gamma _{BB^{}+\phi }=\frac{1}{8\pi }\frac{p_\phi }{MM^{}}\overline{^2}=\frac{1}{8\pi }\frac{p_\phi ^3}{MM^{}}\overline{๐^2}$$
(7)
up to linear order in $`m_s`$. The โbarโover the amplitude squared denotes averaging over initial and summing over final spin (and, if explicitly indicated, over isospin). Anticipating linear momentum dependence of the decay amplitude $``$ we have introduced reduced amplitude $`๐`$ which does not depend on the meson momentum $`p_\phi `$.
Soliton models can be used to calculate the matrix element $``$. Explicitly
$$\Gamma _{BB^{}+\phi }=\frac{3G_{}^2}{8\pi M_BM_B^{}}C_{BB^{}+\phi }^{}p_\phi ^3.$$
For antidecuplet decays ($`=\overline{10}`$):
$$G_{\overline{10}}=G_0G_11/2G_2,C_{\Theta ^+N+K}^{\overline{10}}=1/5,$$
(8)
whereas for decuplet ($`=10`$):
$$G_{10}=G_0+1/2G_2,C_{\Delta N+\pi }^{10}=1/5.$$
(9)
In the nonrelativistic small soliton limit $`^\mathrm{?}`$ in which chiral quark soliton model reproduces many results of the nonrelativistic quark model $`G_1/G_0=4/5`$, $`G_2/G_0=2/5`$ and $`G_{\overline{10}}0`$! This nonintuitive cancellation $`^\mathrm{?}`$ explains the small width of antidecuplet as compared to the one of $`10`$ for example.
One problem concerning this cancellation is that formally
$$G_0๐ช(N_c^{3/2})+๐ช(N_c^{1/2}),G_{1,2}๐ช(N_c^{1/2})$$
(10)
and it looks as if the cancellation were accidental as it occurs between terms of different order in $`N_c`$. For arbitrary $`N_c`$ antidecuplet $`\overline{10}=(0,3)`$ generalizes to $`\mathrm{"}\overline{10}\mathrm{"}=(0,\frac{N_c+3}{2})`$, decuplet $`\mathrm{"}10\mathrm{"}=(3,\frac{N_c3}{2})`$ and octet $`\mathrm{"}8\mathrm{"}=(1,\frac{N_c1}{2})`$ $`^\mathrm{?}`$, and the pertinent Clebsch-Gordan coefficients in fact depend on $`N_c`$:
$$G_{\mathrm{"}\overline{10}\mathrm{"}}=G_0(N_c+1)/4G_11/2G_2.$$
(11)
So the subleading $`G_1`$term is enhanced by additional factor of $`N_c`$ and the cancellation is consistent with $`N_c`$ counting $`^\mathrm{?}`$.
Unfortunately there is another problem concerning the $`N_c`$ counting of the decay width. Because of (6)
$$p_\pi ๐ช(1/N_c),p_K๐ช(1)$$
(12)
and consequently
$$\Gamma _{\Delta N+\pi }๐ช(1/N_c^2),\Gamma _{\Theta ^+N+K}๐ช(1)$$
(13)
in the chiral limit. This $`N_c`$ counting contradicts experimental findings which suggest $`\Gamma _{\Theta ^+N+K}\Gamma _{\Delta N+\pi }`$.
## 4 Closing remarks
Let us finish by summarizing and by adding some remarks.
Experimental situation concerning the existence of exotic baryons is unclear. The new data on photoproduction on deuteron from LEPS with positive evidence have been presented on various conferences but not published. Soon the similar data from G10 experiment at JLab will be released, however the decisive experiment would be certainly โ if ever completed โ high resolution KN scattering experiment.
Little is known about the production mechanism of exotics. Ironically this is an important factor in understanding present experimental situation.
Most models agree that spin of antidecuplet is 1/2 in agreement with $`\chi `$SM prediction. On the contrary parity is a distinguishing feature. Were the parity of $`\mathrm{\Theta }^+`$ positive as in the $`\chi `$SM some other models and some lattice calculations would require revision.
The smallness of the width is very unnatural, although $`\chi `$SM provides formal explanation. If the primary decay coupling of $`\overline{10}8`$ is indeed very small, the SU(3) relations between the decay rates of different members of antidecuplet will not hold due to the flavor representation mixing.
Other members of antidecuplet should be found. This concerns not only $`\mathrm{\Xi }_{\overline{10}}`$ but also five quark cryptoexotic $`\mathrm{\Sigma }`$ and N$``$like states. Recent data on phototoexcitation of nucleon resonances from GRAAl $`^\mathrm{?}`$ may be interpreted as a new narrow antidecuplet N resonance at 1680 MeV. GRAAl sees resonant structure only on neutron but not on proton. This can be understood in terms of magnetic transition moment $`^\mathrm{?}`$ $`\mu _{8\overline{10}}`$ which is proportional to $`Q1`$. Similarly modified PWA $`^\mathrm{?}`$ of $`\pi `$N scattering indicates that such a state migth exist, also STAR data show some structure in the same energy range.
There is no strong theoretical argument against pentaquarks except its unnaturally small width. But โ as recent plethora of theoretical papers shows โ theoretical explanation may be found in many different models. So if high precision experiments will not find $`\mathrm{\Theta }^+`$ and its partners, this may be even more difficult to understand than the small widths and the small mass.
## Acknowledgments
The author would like to thank the organizers of the Rencontres de Moriond for support. The present work was also partially supported by the Polish State Committee for Scientific Research (KBN) under grant 2 P03B 043 24.
## References
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# Untitled Document
An Invariant of Finite Group Actions on Shifts of Finite Type
Daniel S. Silver and Susan G. Williams
The authors were partially supported by NSF grants DMS-0071004 and DMS-0304971. The first author was partly supported by a grant from EGIDE, and the second by a grant from CNRS.2000 Mathematics Subject Classification. Primary 37B10; secondary 20F38, 57M27, 57R56.
Abstract: We describe a pair of invariants for actions of finite groups on shifts of finite type, the left-reduced and right-reduced shifts. The left-reduced shift was first constructed by U. Fiebig, who showed that its zeta function is an invariant, and in fact equal to the zeta function of the quotient dynamical system. We also give conditions for expansivity of the quotient, and applications to combinatorial group theory, knot theory and topological quantum field theory. Keywords: Shift of finite type, knot, representation shift, TQFT.
1. Introduction. Let $`(X,\sigma )`$ be shift of finite type (see section 2 for definitions). By an action of a finite group $`G`$ on $`X`$ we mean a homomorphism $`\varphi `$ of $`G`$ into the automorphism group of $`X`$. We say that $`G`$-actions $`(X,\varphi )`$ and $`(Y,\psi )`$ are conjugate if there is a topological conjugacy $`\eta :XY`$ with $`\eta \varphi (g)=\psi (g)\eta `$ for all $`g`$ in $`G`$.
Identifying points in the same $`G`$-orbit gives a quotient dynamical system $`X/\varphi `$, with the quotient topology and homeomorphism induced by $`\sigma `$. This need not be a shift of finite type, or even an expansive dynamical system. In fact we show in section 4 that for irreducible $`X`$, $`X/\varphi `$ is a shift of finite type if the quotient map is constant-to-one, and nonexpansive otherwise. Ulf-Rainer Fiebig \[Fi93\] showed how to construct a shift of finite type that is an equal-entropy factor of $`X`$ and has the same number of period $`n`$ points, for every $`n`$, as the quotient dynamical system $`X/\varphi `$. In section 3 we examine Fiebigโs shift and a mirror variant, which we call the left-reduced and right-reduced shifts $`{}_{\varphi }{}^{}X`$ and $`X_\varphi `$ of the pair $`(X,\varphi )`$. Fiebigโs construction depends, a priori, on the choice of a presentation of $`X`$ on which the action of $`G`$ is by one-block automorphisms. A satisfying intrinsic definition of the reduced shifts remains elusive. However, we show in Theorem 3.6 that the conjugacy classes of $`{}_{\varphi }{}^{}X`$ and $`X_\varphi `$ are independent of the choice of presentation for $`X`$ and are thus invariants of the group action $`(X,\varphi )`$. Theorem 3.8 and Corollary 3.9 describe the behavior of these invariants under resolving and closing factor maps.
Our work was motivated by questions arising in applications of symbolic dynamics to combinatorial group theory and topology. We describe some of these applications in section 5. In particular, a theorem of Patrick Gilmer \[Gi99\] concerning a class of topological quantum field theories defined by Frank Quinn \[Qu95\] is recovered as a special case of Theorem 3.6.
We thank Mike Boyle, Franรงois Blanchard and especially Bruce Kitchens for valuable discussions. We thank the referee for helpful comments. We are grateful to the University of Maryland, the Institut de Mathรฉmatiques de Luminy and the Centre de Mathรฉmatiques et Informatique, Marseille, for their hospitality while much of this work was being done.
2. Background. In this section we briefly review some basic notions of symbolic dynamics. For more background, includings zeta functions, Bowen-Franks groups and the technique of state-splitting used in the proof of Theorems 3.6 and 3.8, we refer the reader to \[LM95\] or \[Ki98\].
Let $`๐`$ be a finite set or alphabet of symbols. The full shift on $`๐`$ is the dynamical system consisting of the space $`๐^Z`$ with the product topology and the left shift homeomorphism $`\sigma `$ given by $`(\sigma x)_i=(x_{i+1})`$. A subshift $`X`$ is a closed $`\sigma `$-invariant set of some full shift. For $`nN`$ the $`n`$-blocks of $`X`$ are the symbol sequences of length $`n`$ that appear as consecutive entries of some $`xX`$.
Given a finite directed graph $`\mathrm{\Gamma }=(V,E)`$ with adjacency matrix $`A`$, the associated shift of finite type (SFT) $`X_A`$ is the subshift of $`E^Z`$ consisting of bi-infinite sequences $`(x_i)`$ of edges that correspond to paths in $`\mathrm{\Gamma }`$. Vertices of $`\mathrm{\Gamma }`$, which form the index set of $`A`$, are also called states of $`X_A`$. If there is an edge from state $`i`$ to state $`j`$ we say $`j`$ is a follower of $`i`$, and $`i`$ a predecessor of $`j`$. We will always assume that every state has a follower and predecessor; otherwise we can remove that state and its adjacent edges without changing $`X_A`$. If $`\mathrm{\Gamma }`$ has no parallel edges, so that $`A`$ is a zero-one matrix, the edge sequence $`(x_i)`$ is determined by the sequence $`(v_i)`$ of initial vertices of these edges. The SFT $`X_A`$ is irreducible if the graph $`\mathrm{\Gamma }`$ is strongly connected, that is, there is a path from any state to any other state. (This is equivalent to topological transitivity of the dynamical system.)
A homomorphism between dynamical systems $`(X,S)`$ and $`(Y,T)`$ is a continuous, map $`\theta :XY`$ with $`\theta S=T\theta `$. Epimorphisms are also called factor maps and isomorphisms are (topological) conjugacies. If $`\mathrm{\Theta }`$ is a map from the set of $`n`$-blocks of a shift space $`(X,\sigma )`$ to the alphabet of a shift space $`(Y,\sigma )`$ and $`mZ`$, we can define a homomorphism $`\theta :XY`$ by $`\theta ((x_i))=(y_i)`$ where $`y_i=\mathrm{\Theta }(x_{im},\mathrm{},x_{im+n1})`$; $`\eta `$ is an $`n`$-block map with memory $`m`$. For convenience we sometimes use the same symbol for the map on $`n`$-blocks and the map on $`X`$; when we define a map via its action on blocks we will take $`m=0`$. Every homomorphism between subshifts is an $`n`$-block map for some $`n`$.
A one-block map $`\theta :XY`$ is right-resolving if whenever $`ab`$ and $`ac`$ are 2-blocks of $`X`$ with $`\mathrm{\Theta }(b)=\mathrm{\Theta }(c)`$ we have $`b=c`$. If $`X`$ is a shift of finite type described by a graph $`\mathrm{\Gamma }`$ with no parallel edges, this is equivalent to the condition that edges with the same initial vertex have distinct images under $`\mathrm{\Theta }`$. A homomorphism of subshifts is right-closing if it does not identify left-asymptotic points, i.e. points $`(x_i)`$ and $`(x_i^{})`$ with $`x_i=x_i^{}`$ for $`iN`$. Right-resolving maps are right-closing. Right-closing factor maps are bounded-to-one and hence preserve topological entropy. Left-resolving and left-closing maps are defined analogously. One-block maps that are both left and right-resolving are bi-resolving and homomorphisms that are left and right-closing are bi-closing. A factor map between irreducible shifts of finite type is constant-to-one if and only if it is bi-closing \[Na83\].
3. The reduced shift of a group action. In what follows we will consider $`G`$-actions $`(X,\varphi )`$, always assuming that $`X`$ is a shift of finite type and $`G`$ is a finite group. When there is no danger of confustion, we will denote the image of a point $`x`$ under $`\varphi (g)`$ by $`gx`$, and its orbit under $`G`$ by $`Gx`$. Consider the special case of a group action $`(X_A,\varphi )`$ where $`A`$ is an $`n\times n`$ matrix over $`\{0,1\}`$, $`X_A`$ is the associated shift of finite type, and each $`\varphi (g)`$ is a one-block automorphism induced by a permutation of the states of $`X_A`$. We will call such an action a permutation action. We write $`gi`$ for the image of state $`i`$ under the permutation $`\varphi (g)`$, and $`Gi`$ for its orbit.
Proposition 3.1. \[AKM85\] Every finite group action on a SFT is conjugate to a permutation action.
Definition 3.2. Given a permutation action $`(X_A,\varphi )`$, we define the associated right-reduced shift to be the SFT $`X_\varphi `$ given by the matrix $`A_\varphi `$ such that its states are $`G`$-orbits of states of $`X_A`$ and $`A_\varphi (Gi,Gj)=_{kGj}A(i,k).`$ This is well defined since if $`i^{}=gi`$ then
$$\underset{kGj}{}A(i^{},k)=\underset{kGj}{}A(gi,gk)=\underset{kGj}{}A(i,k).$$
Analogously, the left-reduced shift $`{}_{\varphi }{}^{}X`$ is given by the matrix $`{}_{\varphi }{}^{}A`$ with the same state space and $`{}_{\varphi }{}^{}A(Gi,Gj)=_{kGi}A(k,j).`$
The left-reduced shift was constructed (but not named) by U.-R. Fiebig \[Fi93\], who proved the rather surprising result that it has the same zeta function as $`X_A/\varphi `$, or equivalently, the same number of period $`n`$ points for every $`n`$. A permutation action on $`X_A`$ is also a permutation action on the inverse shift given by the transpose $`A^t`$, and it is easy to see that the left-reduced shift for $`(X_{A^t},\varphi )`$ is the inverse of the right-reduced shift for $`(X_A,\varphi )`$. Thus the left-reduced and right-reduced shifts have analogous properties, although Example 3.7 shows that they need not be either conjugate or inverse conjugate. For simplicity, we state and prove most of our results for the right-reduced shift, which we will refer to as the reduced shift when there is no danger of confusion. (This dextrocentric preference follows a tradition in the literature of shifts of finite type, and is also in convenient agreement with the work of Quinn and Gilmer that we discuss in Section 5.)
The construction of the reduced shift, despite its simplicity, is not natural in the sense that it relies on the presentation of the group action as a permutation action. We can easily describe a right-resolving factor map $`\eta :X_AX_\varphi `$: for each pair of states $`i`$, $`j`$ of $`X_A`$, pick a bijection from the set of edges $`(i,k)`$ with $`kGj`$ to the set of edges in $`X_\varphi `$ from $`Gi`$ to $`Gj`$. However, different choices may give nonisomorphic factor maps. (Factor maps $`\eta `$, $`\theta `$ from $`X`$ to $`Y`$ are isomorphic if there are automorphisms $`\alpha `$, $`\beta `$ of $`X`$ and $`Y`$ respectively with $`\theta \alpha =\beta \eta `$.) In general, neither $`X_\varphi `$ nor $`X/\varphi `$ appears in a natural way as a factor of the other.
Example 3.3. Let $`X`$ be the full shift on the elements of the symmetric group $`S_3`$. Let $`G=S_3`$ act on the states of $`X`$ by conjugation, so that $`gG`$ takes the state $`v`$ to $`g^1vg`$. This action induces a permutation action $`\varphi `$ of $`G`$ on $`X`$. The states of the reduced shift $`X_\varphi `$ are the three conjugacy classes of elements of $`S_3`$. If we list these in the order $`[(1)],[(12)],[(123)]`$ then $`X_\varphi `$ is given by the matrix
$$A_\varphi =\left(\begin{array}{ccc}1& 3& 2\\ 1& 3& 2\\ 1& 3& 2\end{array}\right).$$
There are many ways to define a 2-block right-resolving factor map from $`X`$ to $`X_\varphi `$. No choice respects the action of $`G`$. For example, suppose the two-block $`(12),(23)`$ is sent to an edge $`e`$ of $`X_\varphi `$ ($`e`$ must be one of the self-loops on the vertex \[(12)\].) To respect the conjugation by (13) we would need to send the 2-block $`(23),(12)`$ to $`e`$ as well, while to respect the conjugation by (123) we would need to send $`(23),(31)`$ to $`e`$. It is easy to see that different choices can give non-isomrophic factor maps: for example, we may send the fixed points of $`X`$ to distinct fixed points of $`X_\varphi `$, or make some identifications.
Example 3.4. Let
$$A=\left(\begin{array}{ccc}1& 1& 1\\ 1& 1& 0\\ 1& 0& 1\end{array}\right)$$
and let $`\varphi `$ be as in the preceding example. Then $`X_\varphi `$ is given by the graph with symbolic matrix
$$\stackrel{~}{A}_\varphi =\left(\begin{array}{cc}a& b+c\\ d& e\end{array}\right).$$
Let $`q`$ denote the quotient map from $`X`$ to $`X/\varphi `$. There is no right-resolving factor map $`\eta :XX_\varphi `$ satisfying $`\eta =\xi q`$ for some homeomorphism $`\xi :X/\varphi X_\varphi `$, since $`\eta `$ could not identify the left-asymptotic points $`1^{\mathrm{}}2^{\mathrm{}}=(\mathrm{},1,1,2,2,\mathrm{})`$ and $`1^{\mathrm{}}3^{\mathrm{}}`$ which are identified by $`q`$. There is also no factor map $`\eta :XX_\varphi `$ satisfying $`q=\xi \eta `$ for some homomorphism $`\xi :X_\varphi X/\varphi `$. For $`\eta `$ would have to take fixed points $`1^{\mathrm{}}`$, $`2^{\mathrm{}}`$, $`3^{\mathrm{}}`$ to $`a^{\mathrm{}}`$, $`e^{\mathrm{}}`$, $`e^{\mathrm{}}`$ respectively. Since $`\eta `$ is an $`n`$-block map for some $`n`$, it would have to identify $`2^{\mathrm{}}1^n2^{\mathrm{}}`$ and $`2^{\mathrm{}}1^n3^{\mathrm{}}`$, which are not identified by $`q`$.
We may obtain $`A_\varphi `$ from $`A`$ as a matrix product. Suppose $`A`$ has $`n`$ states and $`\overline{n}`$ $`G`$-orbits of states. Let $`V_\varphi `$ be the $`n\times \overline{n}`$ matrix with $`V_\varphi (i,Gi)=1`$ for all $`i`$ and the remaining entries 0. Take $`U_\varphi `$ to be any left inverse of $`V_\varphi `$. Thus $`U_\varphi `$ is $`\overline{n}\times n`$ and $`U_\varphi (Gi,i_0)=1`$ for a selected representative $`i_0`$ of the orbit $`Gi`$. One verifies easily that $`A_\varphi =U_\varphi AV_\varphi `$. Letting $`P_{\varphi (g)}`$ denote the matrix of the permutation $`\varphi (g)`$, we have $`P_{\varphi (g)}V_\varphi =V_\varphi `$ for all $`gG`$.
An elementary strong shift equivalence between square non-negative integral matrices $`A`$ and $`B`$ consists of a pair $`(R,S)`$ of non-negative integral matrices (not necessarily square) with $`RS=A`$ and $`SR=B`$. Matrices $`A`$ and $`B`$ are strong shift equivalent (SSE) if they are linked by a chain of elementary strong shift equivalences, that is, if there are square nonnegative integral matrices $`A=A_0,A_1,\mathrm{}A_n=B`$ with an elementary strong shift equivalence from $`A_{i1}`$ to $`A_i`$ for $`1=1,\mathrm{},n`$. The Decomposition Theorem of R.F. Williams \[Wi73\] states that SFT $`X_A`$ and $`X_B`$ are conjugate if and only if $`A`$ and $`B`$ are SSE.
To an elementary SSE $`(R,S)`$ of 0-1 matrices $`A`$, $`B`$ we may canonically associate a conjugacy from $`X_A`$ to $`X_B`$, as described in \[Ki98\], Lemma 2.1.16, or \[LM95\], p.228. Briefly, since $`A=RS`$, for each nonzero entry $`(i,j)`$ of $`A`$ there is a unique state $`k`$ of $`B`$ with $`R(i,k)=S(k,j)=1`$. This allows us to associate to each $`xX_A`$ a bi-infinite sequence of states of $`B`$. The identity $`SR=B`$ implies that there is a unique edge sequence $`yX_B`$ connecting this sequence of states. The map $`\eta (x)=y`$ is a conjugacy. We observe that this conjugacy will carry a permutation action $`\varphi `$ of $`G`$ on $`X_A`$ to a permutation action $`\psi `$ of $`G`$ on $`X_B`$ if and only if $`RP_{\psi (g)}=P_{\varphi (g)}R`$ and $`SP_{\varphi (g)}=P_{\psi (g)}S`$ for all $`gG`$. For notational simplicity we will view $`X_A`$ and $`X_B`$ as having disjoint state spaces $`I`$ and $`I^{}`$ that each admit an action of $`G`$ by permutations, so that $`gi`$ denotes $`\varphi (g)i`$ for $`iI`$ and $`\psi (g)i`$ for $`iI^{}`$.
Lemma 3.5. Let $`(R,S)`$ be an elementary SSE of 0-1 matrices $`A`$ and $`B`$ that induces a conjugacy of permutation actions $`\varphi `$, $`\psi `$ on $`X_A`$, $`X_B`$ respectively. Then the reduced shifts $`X_\varphi `$ and $`X_\psi `$ are conjugate.
Proof. We claim that the pair $`(U_\varphi RV_\psi ,U_\psi SV_\varphi )`$ is an elementary SSE between $`A_\varphi `$ and $`B_\psi `$. For each $`gG`$,
$$P_{\psi (g)}SV_\varphi =SP_{\varphi (g)}V_\varphi =SV_\varphi ,$$
which implies that the $`i`$-th and $`gi`$-th rows of $`SV_\varphi `$ are identical for all $`i`$. Now, $`V_\psi U_\psi `$ has 1 in the $`(i,i_0)`$ entry where $`i_0`$ is the distinguished representative of the orbit $`Gi`$, and 0 elsewhere. Hence $`V_\psi U_\psi SV_\varphi =SV_\varphi `$. This yields
$$(U_\varphi RV_\psi )(U_\psi SV_\varphi )=U_\varphi RSV_\varphi =U_\varphi AV_\varphi =A_\varphi .$$
Similarly, $`(U_\psi SV_\varphi )(U_\varphi RV_\psi )=B_\psi `$.
Theorem 3.6. If $`(X_A,\varphi )`$ and $`(X_B,\psi )`$ are conjugate permutation actions on SFT then the reduced shifts $`X_\varphi `$ and $`X_\psi `$ are conjugate.
Proof. By the Decomposition Theorem, the conjugacy can be expressed as a composition of conjugacies corresponding to elementary SSE. In light of the preceding lemma, it suffices to show that the decomposition can be carried out in such a way that each elementary SSE induces a conjugacy of permutation actions.
We follow the proof of the Decomposition Theorem in \[LM95\] (Theorem 7.1.2). By passing from $`X_A`$ by a chain of elementary SSE to a higher block presentation we may replace the original conjugacy by a 1-block map. The $`n`$-block presentation of $`X_A`$ is again given by a 0-1 matrix, and each permutation $`\varphi (g)`$ of the states of $`X_A`$ naturally induces a permutation of $`n`$-blocks that is a conjugate permutation action on the $`n`$-block presentation.
Thus we may assume we have a 1-block conjugacy $`\eta :X_AX_B`$. If $`\eta ^1`$ is also a 1-block map then the conjugacy is simply a renaming of symbols, and the result is clear. If $`\eta ^1`$ has anticipation $`k>1`$ we use an out-splitting of the graph of $`X_A`$ to reduce the anticipation. At each vertex of the graph of $`X_A`$ we partition the outgoing edges according to their images under $`\eta `$. These partition elements are the states of a conjugate shift $`X_{\stackrel{~}{A}}`$, where $`\stackrel{~}{A}`$ is a 0-1 matrix, and $`\eta `$ induces a one-block conjugacy $`\stackrel{~}{\eta }`$ from $`X_{\stackrel{~}{A}}`$ to the two-block presentation of $`X_B`$ such that $`\stackrel{~}{\eta }^1`$ has the same memory as $`\eta `$ but anticipation $`k1`$. Since $`\eta `$ intertwines the $`G`$-actions on $`X_A`$ and $`X_B`$, if two edges leaving state $`i`$ have the same image under $`\eta `$ then any $`gG`$ carries them to edges leaving state $`gi`$ with the same image under $`\eta `$. This determines a permutation action on $`X_{\stackrel{~}{A}}`$. It is easy to see that all of these conjugacies preserve the $`G`$-actions.
The memory of $`\eta ^1`$ may be reduced by in-splittings in an analogous fashion. An induction argument finishes the proof.
In view of Proposition 3.1 and Theorem 3.6, we can speak of the (left or right) reduced shift of an arbitrary $`G`$-action on a SFT with the understanding that it is well defined up to topological conjugacy.
We say the $`G`$-action $`(Y,\psi )`$ is a factor of the $`G`$-action $`(X,\varphi )`$ if there is a factor map $`\eta :XY`$ with $`\eta \varphi (g)=\psi (g)\eta `$ for all $`gG`$. In this case the quotient dynamical system $`Y/\psi `$ is a factor of $`X/\varphi `$, and it is natural to ask whether the reduced shift $`Y_\psi `$ is a factor of $`X_\varphi `$. Example 3.7 shows that this can fail even for almost invertible factor maps between irreducible SFT. Theorem 3.8 and its corollary give a positive answer for right-resolving and right-closing factor maps. Analogous results hold for left-closing and left-resolving factors.
Example 3.7. Let
$$A=\left(\begin{array}{cccccc}1& 0& 1& 0& 1& 0\\ 0& 1& 0& 1& 0& 1\\ 1& 1& 1& 0& 0& 0\\ 1& 1& 0& 1& 0& 0\\ 1& 1& 0& 0& 1& 0\\ 1& 1& 0& 0& 0& 1\end{array}\right)$$
and let $`\varphi `$ be the permutation action on $`X_A`$ of the cyclic group $`GZ/4`$ generated by the permutation $`(12)(3456)`$. Then
$$A_\varphi =\left(\begin{array}{cc}1& 2\\ 2& 1\end{array}\right),_\varphi A=\left(\begin{array}{cc}1& 1\\ 4& 1\end{array}\right).$$
Here we have taken the states of the reduced shifts to be $`G1=\{1,2\}`$ and $`G3=\{3,4,5,6\}`$ in that order. The left-reduced shift is conjugate to neither the right-reduced shift nor its inverse, as they have non-isomorphic Bowen-Franks groups $`Z^2/(IA_\varphi )Z^2Z/2Z/2`$ and $`Z^2/(I_\varphi A)Z^2Z/4`$.
Let
$$B=\left(\begin{array}{ccccc}1& 1& 1& 1& 1\\ 1& 1& 0& 0& 0\\ 1& 0& 1& 0& 0\\ 1& 0& 0& 1& 0\\ 1& 0& 0& 0& 1\end{array}\right).$$
A one-block factor map of $`X_A`$ onto $`X_B`$ is obtained by identifying the first two states of $`A`$, and this induces a $`Z/4`$-action $`\psi `$ on $`X_B`$ that cyclically permutes the last four states of $`B`$. Points of $`X_B`$ that are right-asymptotic to the fixed point on the first state have two preimages in $`X_A`$ but every other point has a unique preimage. We have
$$B_\psi =\left(\begin{array}{cc}1& 4\\ 1& 1\end{array}\right).$$
The right-reduced shift $`X_\psi `$ is not a factor of $`X_\varphi `$, as can be seen from either Theorem 4.2.16 or Theorem 4.2.19 of \[Ki98\]. (The latter theorem says that a factor map between irreducible SFT induces an epimorphism of their Bowen-Franks groups.)
A simpler, but reducible, example of these phenomena is given by
$$A=\left(\begin{array}{ccc}1& 0& 1\\ 0& 1& 1\\ 0& 0& 1\end{array}\right)$$
with permutation action of $`GZ/2`$ generated by the permutation $`(12)`$. Then
$$A_\varphi =\left(\begin{array}{cc}1& 1\\ 0& 1\end{array}\right),_\varphi A=\left(\begin{array}{cc}1& 2\\ 0& 1\end{array}\right).$$
Let $`B`$ be the quotient $`{}_{\varphi }{}^{}A`$ with trivial $`G`$ action $`\psi `$. Then $`A_\varphi `$ and $`{}_{\varphi }{}^{}A`$ are nonconjugate, and $`B_\psi =B`$ is not a quotient of $`A_\varphi `$.
Theorem 3.8. Suppose the permutation action $`(X_B,\psi )`$ is a factor of the permutation action $`(X_A,\varphi )`$ by a right-resolving 1-block map $`\eta `$. Then there are right resolving 1-block maps $`\overline{\eta }:X_\varphi X_\psi `$, $`\theta _1:X_AX_\varphi `$ and $`\theta _2:X_BX_\psi `$ with $`\theta _2\eta =\overline{\eta }\theta _1`$.
Proof. The 1-block map $`\eta `$ induces a graph homomorphism from the graph of $`X_A`$ to the graph of $`X_B`$, which takes $`i`$ to a vertex we denote by $`\eta (i)`$. Since $`\eta `$ is right-resolving, distinct edges with the same initial vertex $`i`$ have distinct images under this graph homomorphism. Since $`\eta `$ preserves the group action, the images of vertices in an orbit $`Gj`$ comprise the orbit $`G\eta (j)`$. Hence $`\eta `$ gives a bijection from the set of edges with initial vertex $`i`$ and terminal vertex in $`Gj`$ to the set of edges with initial vertex $`\eta (i)`$ and terminal vertex in $`G\eta (j)`$. This means the number of edges in the graph of $`X_\varphi `$ from $`Gi`$ to $`Gj`$ is equal to the number of edges in the graph of $`X_\psi `$ from $`G\eta (i)`$ to $`G\eta (j)`$; we define the 1-block map $`\overline{\eta }`$ by making any choice of bijections. We may also let $`\theta _2`$ be given by any choice of bijections from the edges in the graph of $`X_B`$ with initial state $`i^{}`$ and terminal state in $`Gj^{}`$ to the edges in the graph of $`X_\psi `$ from state $`Gi^{}`$ to state $`Gj^{}`$. There is now a unique 1-block map $`\theta _1`$ which gives the desired commutativity.
Because $`\eta `$ and $`\theta _2`$ are right-resolving, the composite map $`\theta _2\eta =\overline{\eta }\theta _1`$ is also right-resolving. Now since $`\theta _1`$ is onto, it is easy to see that $`\overline{\eta }`$ must be right-resolving as well.
Corollary 3.9. If the $`G`$-action $`(Y,\psi )`$ is a factor of the $`G`$-action $`(X,\varphi )`$ by a right-closing map $`\eta `$ then $`Y_\psi `$ is a right-closing factor of $`X_\varphi `$.
Proof. By Proposition 3.1, we can assume from the start that the actions are permutation actions. By Proposition 1 of \[BKM85\], there is a topological conjugacy $`\pi :\stackrel{~}{X}X`$ such that $`\eta \pi `$ is a right-resolving 1-block map. It is easy to see from their construction that $`\pi ^1`$ induces a permutation action on $`\stackrel{~}{X}`$. We apply Theorem 3.8 to $`\eta \pi `$, then use Theorem 3.6 together with the fact that the composition of a right-resolving map with a topological conjugacy is right-closing.
4. $`G`$-stabilizers and $`G`$-orbits of periodic points. Given a $`G`$-action $`(X,\varphi )`$, for $`xX`$ we denote by $`\mathrm{Stab}(x)`$ the $`G`$-stabilizer of $`x`$, that is, the subgroup of all $`gG`$ with $`gx=x`$. For permutation actions, the $`G`$-stabilizer of a state or word of $`X`$ may be defined similarly. In this case the $`G`$-stabilizer of $`x=(x_i)`$ is $`\mathrm{Stab}(x)=_{iZ}\mathrm{Stab}(x_i)`$. The number of preimages of a point $`GxX/\varphi `$ under the quotient map is the index in $`G`$ of $`\mathrm{Stab}(x)`$.
Theorem 4.1. Let $`(X,\varphi )`$ be a $`G`$-action on an irreducible shift of finite type. (i) If every $`xX`$ has the same stabilizer then the quotient map is constant-to-one, and $`X/\varphi `$, $`X_\varphi `$ and $`{}_{\varphi }{}^{}X`$ are all conjugate shifts of finite type. (ii) If some pair of points of $`X`$ have different $`G`$-stabilizers then $`X/\varphi `$ is nonexpansive.
Proof. We may assume $`(X,\varphi )`$ is a permutation action. Suppose first that every point has the same $`G`$-stabilizer $`H`$. Since every state of $`X`$ appears in some periodic point of period at most $`n`$, the number of states of $`X`$, by passing to the $`n`$-block presentation of $`X`$ we can assume every state of $`X`$ has stabilizer $`H`$. Now it is clear that the one-block map $`iGi`$ is a bi-resolving map from $`X`$ to $`{}_{\varphi }{}^{}X=X_\varphi `$, and the image is topologically conjugate to $`X/\varphi `$.
Now suppose there are points of $`X`$ with different stabilizers. We first show there must be periodic points $`u^{\mathrm{}}`$, $`v^{\mathrm{}}`$ with different stabilizers. Since $`X`$ is irreducible there is a periodic point $`v^{\mathrm{}}`$ containing all states, so that $`\mathrm{Stab}(v^{\mathrm{}})`$ is a subgroup of $`\mathrm{Stab}(x)`$ for all $`xX`$. Choose any $`yX`$ with $`\mathrm{Stab}(y)\mathrm{Stab}(v^{\mathrm{}})`$. Then $`y`$ contains a word $`u`$ such that $`u^{\mathrm{}}X`$, and $`\mathrm{Stab}(y)`$ is a subgroup of $`\mathrm{Stab}(u^{\mathrm{}})`$.
Let $`g\mathrm{Stab}(u)\mathrm{Stab}(v)`$. We can find words $`w,w^{}`$ such that $`uwv`$ and $`vw^{}u`$ are words of $`X`$. Then (gu)(gw)(gv)=u(gw)(gv) is also a word of $`X`$. For each positive integer $`m`$ set
$$\begin{array}{cc}\hfill x^{(m)}& =v^{\mathrm{}}w^{}.u^{2m+1}wv^{\mathrm{}}\hfill \\ \hfill y^{(m)}& =v^{\mathrm{}}w^{}.u^{2m+1}(gw)(gv)^{\mathrm{}}.\hfill \end{array}$$
(Here the point precedes the 0-coordinate.) Then $`x^{(m)}`$ and $`y^{(m)}`$ have different images $`[x^{(m)}]`$, $`[y^{(m)}]`$ under the quotient map. However, for every $`nZ`$ the central $`(2m+1)`$-block of $`\sigma ^ny^{(m)}`$ agrees with the central $`(2m+1)`$-block of either $`\sigma ^nx^{(m)}`$ or $`\sigma ^ngx^{(m)}`$, so that $`[\sigma ^nx^{(m)}]`$ and $`[\sigma ^ny^{(m)}]`$ are close in the quotient topology. Hence there is no expansive constant for $`X/\varphi `$.
We next give a formula for the number of $`G`$-orbits of period $`n`$ points of $`X`$. Note that this is different from the number of period $`n`$ points of the quotient dynamical system $`X/\varphi `$, since the quotient map may change the period of a point. An application of this result appears as Propositon 5.1 below.
For each $`gG`$ the set Fix$`(g)`$ of points of $`X`$ fixed by $`g`$ is again a SFT. If we assume that $`\varphi `$ is a permutation action on $`X_A`$ then Fix$`(g)=X_{A_g}`$ where $`A_g`$ is the principal submatrix of $`A`$ corresponding to the set of symbols fixed by $`g`$. (If this set is empty we take $`A_g=(0)`$.)
Proposition 4.2. Let $`G`$ act by permutations on a shift of finite type $`X_A`$. The number $`N_n`$ of $`G`$-orbits of period $`n`$ points of $`X_A`$ is given by
$$N_n=\frac{1}{|G|}\underset{gG}{}\mathrm{trace}(A_g^n).$$
Hence the sequence $`\{N_n\}`$ satisfies a linear homogeneous recurrence relation with constant coefficients.
Proof. A combinatorial result of Cauchy and Frobenius commonly known as the Burnside Lemma (cf. \[DM96\], p.24) says that if a group $`G`$ acts on a set $`S`$ then the number of $`G`$-orbits of $`S`$ is
$$\frac{1}{|G|}\underset{gG}{}|\mathrm{Fix}(g)|,$$
where $`\mathrm{Fix}(g)`$ is the set of points fixed by $`g`$. If we take $`S`$ to be the set of period $`n`$ points of $`X_A`$ then $`\mathrm{Fix}(g)`$ is the set of period $`n`$ points of $`X_{A_g}`$, which has cardinality trace$`(A_g^n)`$. The sequence $`\{\mathrm{trace}(A_g^n)\}`$ satisfies the linear recurrence with characteristic polynomial $`det(ItA_g)`$, so the sum satisfies the recurrence relation given by the least common multiple of these polynomials.
5. Representation shifts, knots and topological quantum field theories. In this section we describe a class of SFT called representation shifts that admit a natural group action, and outline applications to knot theory and topological quantum field theory.
Assume that $`\mathrm{\Pi }`$ is a finitely presented group with epimorphism $`\chi :\mathrm{\Pi }Z`$, and let $`x\chi ^1(1)`$. Then $`\mathrm{\Pi }`$ can be described as an HNN extension $`x,Bx^1ax=\varphi (a),aU`$. Here $`B`$ is a finitely generated subgroup of $`K=\mathrm{ker}\chi `$, and the map $`\varphi `$ is an isomorphism between finitely generated subgroups $`U,V`$ of $`B`$. The subgroup $`B`$ is an HNN base, $`x`$ is a stable letter, $`\varphi `$ is an amalgamating map. Details can be found in \[LS77\].
Conjugation by $`x`$ induces an automorphism of $`K`$. Letting $`B_j=x^jBx^j,U_j=x^jUx^j`$ and $`V_j=x^jVx^j,jZ`$, we can express $`K`$ as an infinite amalgamated free product
$$K=B_jV_j=U_{j+1},jZ.$$
For any finite group $`G`$, the set $`\mathrm{Hom}(K,G)`$ of representations of $`K`$ in $`G`$ may be viewed as a SFT. The state set is $`\mathrm{Hom}(U,G)`$; an edge is an element $`\rho _0\mathrm{Hom}(B,G)`$ with initial state $`\rho _0|_U`$ and terminal state $`\rho _0|_V\varphi `$. (Note that the cardinality of the edge set is bounded by $`|G|^m`$ where $`m`$ is the cardinality of a generating set of $`B`$.) If $`\rho =(\rho _j)`$ is a bi-infinite path then the representations from $`B_j`$ to $`G`$ given by $`y\rho _j(x^jyx^j)`$ have a unique common extension to an element of $`\mathrm{Hom}(K,G)`$ that we will also denote by $`\rho `$.
We call this SFT the representation shift of $`K`$ in $`G`$ and denote it by $`\mathrm{\Phi }_G=\mathrm{\Phi }_G(\mathrm{\Pi },\chi ,x)`$ (see \[SW96\]). The usual topology on the SFT coincides with the compact-open topology on $`\mathrm{Hom}(K,G)`$, and the shift map $`G`$ can be described by $`\sigma \rho (y)=x^1yx`$. It is clear from this intrinsic description that the topological conjugacy class of $`\mathrm{\Phi }_G`$ is independent of the choice of HNN base $`B`$.
The group $`G`$ acts on $`\mathrm{\Phi }_G`$ by inner automorphism of the image space, $`(g\rho )(x)=g^1\rho (x)g`$. The case where $`G`$ is the symmetric group $`S_n`$ is of particular interest in studying finite-index subgroups of $`K`$ (see \[SW96\], \[SW99\] and \[SW04\]). In this case the state space of the representation shift typically grows very quickly with $`n`$. Since the group of inner automorphisms is isomorphic to $`S_n`$, the corresponding reduced shift will be considerably simpler. (For $`n6`$ all automorphisms of $`S_n`$ are inner: see theorem 6.20 of \[Is94\]).
Much of our original motivation came from the study of knots. A knot $`k`$ is a smoothly embedded circle in the $`3`$-sphere $`S^3`$. For convenience, we assume that $`k`$ is oriented. Two knots are regarded as the same if they are ambiently isotopic. Although the knot group $`\mathrm{\Pi }=\pi _1(S^3k)`$ is essentially a complete invariant, unlocking all of its information is not a reachable task at present. Fortunately, many tractable invariants can be computed from the group.
Let $`x\mathrm{\Pi }`$ be the element represented by a meridian curve encircling $`k`$ with linking number $`1`$. The abelianization of $`\mathrm{\Pi }`$ is infinite cyclic, and we let $`\chi :\mathrm{\Pi }Z`$ be the abelianization homomorphism that maps $`x`$ to $`1`$. The kernel $`K`$ of $`\chi `$ is the commutator subgroup of $`\mathrm{\Pi }`$. Given a finite group $`G`$, the representation shift $`\mathrm{\Phi }_G=\mathrm{\Phi }_G(\mathrm{\Pi },\chi ,x)`$ is an invariant of $`k`$ \[SW96\] (see also \[SW99\]).
For any positive integer $`n`$, the set $`\mathrm{Fix}(\sigma ^n)`$ of period $`n`$ points coincides with the set $`\mathrm{Hom}(\pi _1M_n,G)`$, where $`M_n`$ is the $`n`$-fold cyclic cover of $`S^3`$ branched over $`k`$. Details can be found in \[SW99\]. Invariants of branched covers $`M_n`$ are invariants of $`k`$. (Such invariants were first considered by J. Alexander and G. Briggs in the early 1900โs, and they remain an important class.) The group $`G`$ acts on $`\mathrm{Hom}(\pi _1M_n,G)`$ by inner automorphism as above, and the orbit set $`\mathrm{Hom}(\pi _1M_n,G)/G`$ can be identified with the set of flat $`G`$-bundles on $`M_n`$. Proposition 4.2 immediately yields the following.
Proposition 5.1. For any knot $`k`$ and finite group $`G`$, the number of flat $`G`$-bundles over the $`r`$-fold cyclic cover $`M_n`$ of $`S^3`$ satisfies a linear recurrence.
Proposition 5.1 should be compared to Theorem 4.2 of \[SW99\] or Proposition 2.3 of \[Gi99\], either of which shows that $`|\mathrm{Hom}(\pi _1M_n,G)|`$ satisfies a linear recurrence.
Gilmerโs paper \[Gi99\] is concerned with topological quantum field theories (TQFT), and was another source of motivation for us. TQFT, which arose from quantum physics, offer a framework for the understanding invariants such as the Jones polynomial for links or Donaldson invariants of $`4`$-manifolds or the discovery of new invariants.
Roughly speaking, a (d+1)-dimensional TQFT assigns to a $`d`$-dimensional oriented manifold $`Y`$, called a space, a module $`Z(Y)`$ over a coefficient ring $`R`$. When spaces $`Y_1`$ and $`Y_2`$ are the โincomingโ and โoutgoingโ boundaries of a spacetime, an oriented $`(d+1)`$-dimensional manifold $`X`$, a homomorphism $`Z_X:Z(Y_1)Z(Y_2)`$ is assigned. In particular, we require $`Z_{Y\times [0,1]}=\mathrm{id}:Z(Y)Z(Y)`$, which implies that if $`Z_X`$ is nontrivial, then $`X`$ is topologically nontrivial (that is, not a product of a $`d`$-dimensional manifold with the unit interval). Various other axioms are imposed. For example, $`Z(Y_1Y_2)=Z(Y_1)_RZ(Y_2)`$, where $``$ denotes disjoint union. Also, if $`X_1`$ has incoming (resp. outgoing) boundaries $`Y_1`$ (resp. $`Y_2`$) while $`X_2`$ has incoming (resp. outgoing) boundaries $`Y_2`$ (resp. $`Y_3`$):
$$Y_1\stackrel{X_1}{}Y_2\stackrel{X_2}{}Y_3,$$
then the associated homomorphisms compose in a natural way:
$$Z_{X_1_{Y_2}X_2}=Z_{X_2}Z_{X_1}.$$
For a more complete discussion, the reader might consult \[Qu95\] or \[At89\]. In Quinnโs very general approach, manifolds can be replaced by finite CW complexes with variously defined boundaries.
Quinn \[Qu95\] uses a finite group $`G`$ to construct a TQFT. For the sake of simplicity, we describe his TQFT only for a special case that arises in knot theory. As above, let $`k`$ be an oriented knot with group $`\mathrm{\Pi }=\pi _1(S^3\mathrm{})`$ and abelianization homomorphism $`\chi :\mathrm{\Pi }Z`$ mapping the distinguished element $`x\mathrm{\Pi }`$ to $`1`$. It is well known that $`S^3k`$ admits a smooth map $`f`$ to $`S^1`$ inducing $`\chi `$ on first homology groups and such that the preimage of a regular value is a connected orientable surface, the interior of a surface $`YS^3`$ with boundary $`k`$, called a Seifert surface for the knot. Cutting $`S^3`$ along $`Y`$ produces a compact manifold $`X`$ with two boundary components $`Y_1,Y_2`$ that are copies of $`Y`$.
Quinnโs TQFT assigns a $`Q`$-vector space $`Z(Y_i)`$ to $`Y_i`$, $`i=1,2`$, with basis consisting of homomorphisms from $`\pi _1(Y_i)`$ to $`G`$ modulo inner automorphisms of $`G`$; in other words, $`Z(Y_i)`$ has basis $`\mathrm{Hom}(\pi _1(Y_i),G)/G`$. Choosing basepoints $`y_iY_i`$ and a path $`s`$ in $`X`$ connecting $`y_1`$ and $`y_2`$, one defines a homomorphism $`s_{}:\pi _1(Y_2,y_2)\pi _1(X,y_1)`$ sending a loop $`\gamma `$ at $`y_2`$ to $`s\gamma s^1`$. For any $`\beta :\pi _1(X,y_1)G`$, let $`\beta _1:\pi _1(Y_1,y_1)G`$ be the composition of the map $`\pi _1(Y_1,y_1)\pi _1(X,y_1)`$ induced by inclusion, and $`\beta `$. Let $`\beta _2:\pi _1(Y_2,y_2)G`$ be the composition of $`s_{}`$ and $`\beta `$. Combining these ingredients, we define $`Z_{Y_1,Y_2}`$ to be the homomorphism from $`Z(Y_1)`$ to $`Z(Y_2)`$ that maps $`[\alpha ]\mathrm{Hom}(\pi _1(Y_1),G)/G`$ to $`_\beta [\beta _2],`$ where the sum is taken over all $`\beta :\pi _1(X,y_1)G`$ such that $`\beta _1=\alpha .`$ In a sense, $`Z_{Y_1,Y_2}`$ records the various ways that $`\beta _2`$ extends over $`\pi _1(X)`$ modulo inner automorphisms of $`G`$.
We identify $`Z_{Y_1,Y_2}`$ with its matrix representation (with respect to the given bases). Since its entries are nonnegative integers, it defines a shift of finite type. In fact, it is the reduced shift of the representation shift $`\mathrm{\Phi }_G=\mathrm{\Phi }_G(\mathrm{\Pi },\chi ,x)`$, defined above. Since the representation shift is well defined, it follows from Theorem 3.6 that the strong shift equivalence class of $`Z_{Y_1,Y_2}`$ is an invariant of the knot $`k`$. This fact was established by Gilmer in \[Gi99\] using topological methods.
References
\[AKM85\] R.L. Adler, B. Kitchens and B.H. Marcus, Finite group actions on shifts of finite type, Ergod. Th. & Dynam. Sys. 5 (1985), 1โ25.
\[At89\] M.F. Atiyah, Topological Quantum Field Theories, Publ. Math. Inst. Hautes Etudes Sci. 68 (1989), 175โ186.
\[BKM85\] M. Boyle, B. Kitchens and B. Marcus, A note on minimal covers for sofic systems, Proc. Amer. Math. Soc. 95 (1985), 403โ411.
\[DM96\] J.D. Dixon and B. Mortimer, Permutation Groups, Springer-Verlag, 1996.
\[Fi93\] U. Fiebig, Periodic points and finite group actions on shifts of finite type, Ergod. Th. & Dynam. Sys. 13 (1993), 485โ514.
\[Gi99\] P.M. Gilmer, Topological quantum field theory and strong shift equivalence, Canad. Math. Bull. 42 (1999), 190โ197.
\[Is94\] I.M. Isaacs, Algebra: A Graduate Course, Brooks/Cole, Belmont, CA, 1994.
\[Ki98\] B.P. Kitchens, Symbolic Dynamics: One-sided, two-sided and countable state Markov chains, Springer-Verlag, Berlin 1998.
\[LM95\] D. Lind and B. Marcus, An Introduction to Symbolic Dynamics and Coding, Cambridge University Press, 1995.
\[LS77\] R.C. Lyndon and P.E. Schupp, Combinatorial Group Theory, Springer-Verlag, Berlin, 1977.
\[Na83\] M. Nasu, Constant-to-one and onto global maps of homomorphisms between strongly connected graphs, Ergod. Th.& Dynam. Sys. 3 (1983), 387-411.
\[Qu95\] F. Quinn, Lectures on axiomatic topological quantum field theory, Geometry and Quantum Field Theory (D, Freed, K. Uhlenbeck, ed.), Amer. Math. Soc., 1995.
\[Si93\] D.S. Silver, Augmented group systems and $`n`$-knots, Mathematische Annalen 296 (1993), 585โ593.
\[SW96\] D.S. Silver and S.G. Williams, Augmented group systems and shifts of finite type, Israel J. Math. 95 (1996), 231โ251.
\[SW99\] D.S. Silver and S.G. Williams, Knot invariants from symbolic dynamical systems, Trans. Amer. Math. Soc. 351 (1999), 3243โ3265.
\[SW99\] D.S. Silver and S.G. Williams, On groups with uncountably many subgroups of finite index, J. Pure Appl. Algebra 140 (1999), 75โ86.
\[SW04\] D.S. Silver and S.G. Williams, Lifting representations of $`Z`$-groups, preprint, April 2004.
\[St51\] N. Steenrod, The Topology of Fibre Bundles, Princeton Univ. Press, Princeton NJ, 1951.
\[Wi73\] R.F. Williams, Classification of subshifts of finite type, Annals of Mathematics 98 (1973), 120โ153; erratum, Annals of Mathematics 99 (1974), 380โ381.
Dept. of Mathematics and Statistics, Univ. of South Alabama, Mobile, AL 36688
email: silver@jaguar1.usouthal.edu, swilliam@jaguar1.usouthal.edu
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# Fluctuation-dissipation relations outside the linear response regime in a two-dimensional driven lattice gas along the direction transverse to the driving force
## Model:
Let $`\eta _i`$ be an occupation variable defined on each site $`i=(i_x,i_y)`$ of a two-dimensional square lattice $`\{(i_x,i_y)|0i_xL,0i_yL\}`$. The variable $`\eta _i`$ is 1 if the $`i`$th site is occupied by a particle, and 0 if it is unoccupied. Periodic boundary conditions are imposed by setting $`\eta _i=\eta _j`$, where $`j=(L,i_y)`$ in the case $`i_x=0`$, and $`\eta _i=\eta _j`$, where $`j=(i_x,L)`$ in the case $`i_y=0`$. The array of all occupation variables, $`\{\eta _i\}`$, is denoted $`๐ผ`$ and called the โconfigurationโ.
The time evolution of $`๐ผ`$ is described by the following rule: At each time step, randomly choose a nearest-neighbor pair $`i,j`$, and exchange the values of $`\eta _i`$ and $`\eta _j`$ with the probability $`c(i,j;๐ผ)=\{1+\mathrm{exp}[\beta Q(๐ผ๐ผ^{ij})]\}^1`$, where $`๐ผ^{ij}`$ is the configuration obtained from $`๐ผ`$ through this exchange, and $`\beta =1/T`$ is the inverse temperature with the Boltzmann constant set to unity. $`Q(๐ผ๐ผ^{ij})`$ represents the heat absorbed from the heat bath as a result of the configuration change $`๐ผ๐ผ^{ij}`$. The total particle number, $`N=_i\eta _i`$, is conserved throughout the time evolution. The density $`\rho =N/L^2`$ is a parameter of the model. Hereafter, we regard the unit of time to be $`L^2`$ time steps, which is the number of time steps for which an arbitrary site is chosen once on average. We refer to this time as 1 MCS (Monte Carlo step per site).
In this paper, we study a two-dimensional DLG with
$$Q(๐ผ๐ผ^{})H_0(๐ผ^{})H_0(๐ผ)Ej_\mathrm{p}(๐ผ๐ผ^{}),$$
(1)
where $`E`$ is an external driving force, and $`H_0(๐ผ)`$ describes an interaction between particles, written $`H_0(๐ผ)_{i,j}\eta _i\eta _j`$, where $`i,j`$ denotes a nearest-neighbor pair. The quantity $`j_\mathrm{p}(๐ผ๐ผ^{})`$ is the spatially-averaged current, that is, the net number of particles flowing in the $`x`$ direction: $`j_\mathrm{p}(๐ผ๐ผ^{})_i[\eta _i(1\eta _i^{})\eta _{i+(1,0)}^{}(1\eta _{i+(1,0)})\eta _i^{}(1\eta _i)\eta _{i+(1,0)}(1\eta _{i+(1,0)}^{})]`$. In this study, we fix $`\beta =0.5`$ in order for the system to be far from the critical region (Note that the critical temperature of the model with $`E=0`$ is $`\beta _\mathrm{c}=1.76`$.), and choose large values of $`E`$ in order for the system to be far from equilibrium.
## Our aim:
In DLGs, regarded as one of the simplest classes of nonequilibrium models, statistical properties of NESS have been investigated from various points of view. Among them, there is an interesting report that a large deviation functional of density fluctuations is shape dependent eyink . In a two-dimensional DLG, although the general properties of fluctuations are quite different from those of equilibrium states, it was found numerically that a fluctuation relation holds HSI , where we consider only properties along the direction transverse to the external force $`E`$. This fluctuation relation is a relation among density fluctuations, the chemical potential chemi ; sst , and the temperature of the environment.
In equilibrium cases, the fluctuation relation is closely related to fluctuation-dissipation relations, which relate dynamical properties of equilibrium fluctuations with transport properties in the linear response regime HSVI . Then noting that the fluctuation relation holds in the DLG even far from equilibrium HSI , we wish to also investigate the validity of fluctuation-dissipation relations far from equilibrium, and determine if their equilibrium forms hold here as well.
In order to obtain such relations, we directly measure transport coefficients and dynamical fluctuations in the direction transverse to the driving force $`E`$ in the two-dimensional DLG investigated above. In spite of the fact that these measured values differ from those for the equilibrium state($`E=0`$), we numerically find that three fluctuation-dissipation relations, the Einstein relation, the fluctuation-response relation, and the Green-Kubo relation, seem to be valid even for NESS far from equilibrium.
## Einstein relation:
In the linear response regime near $`E=0`$, the Einstein relation for interacting many-body systems is written
$$D\chi =\sigma T,$$
(2)
where $`D`$ is the density diffusion constant, $`\chi `$ is the intensity of density fluctuations, and $`\sigma `$ is the conductivity kubo .
In order to investigate the validity of (2) in the direction transverse to the external driving force $`E`$ (the $`y`$-direction), we first need to define the density diffusion constant $`D`$ as the coefficient of the diffusion term in the evolution equation describing the averaged behavior of a density field in this direction. In this paper, we consider the case in which we have prepared as the initial state, a steady state under the perturbation potential
$$V(i_y)=\mathrm{\Delta }\mathrm{sin}\frac{2\pi i_y}{L}$$
(3)
which is obtained by adding $`_i\eta _iV(i_y)`$ to $`H_0(๐ผ)`$. Then, we remove $`V(i_y)`$ at $`t=0`$ in order to measure the relaxation of the density field $`\widehat{\rho }(t)`$. The function $`\widehat{\rho }(t)`$ is the Fourier transform of the coarse-grained density $`\rho (i_y)`$. These quantities are defined as $`\widehat{\rho }(t)_{i_y=1}^L\rho (i_y)\mathrm{sin}\frac{2\pi i_y}{L}`$, and $`\rho (i_y)\frac{1}{L}_{i\mathrm{\Omega }_{i_y}}\eta _i`$, where $`\mathrm{\Omega }_{i_y}=\{(i_x,i_y)|1i_xL,i_y\}`$.
In Fig. 1, choosing $`\mathrm{\Delta }`$ as a sufficiently small value, $`\mathrm{ln}(\widehat{\rho }(t)_E^V/\mathrm{\Delta })`$ is plotted as a function of $`t`$ in the case $`(\rho ,E)=(0.5,10)`$ with $`L=32`$ and $`\mathrm{\Delta }=0.2`$. Here, $`_E^V`$ represents the statistical average under the relaxation process. Because exponentially decaying behavior of $`\widehat{\rho }(t)_E^V`$ is observed, $`D`$ can be estimated from the form
$$\widehat{\rho }(t)_E^V=\mathrm{const}.\mathrm{e}^{D\left(\frac{2\pi }{L}\right)^2t}.$$
(4)
In the inset of Fig. 1, $`D`$ is plotted as a function of the system size $`L`$ in the cases $`E=0`$ and $`E=10`$, with $`\rho =0.5`$. Because both values of $`D`$ seem to converge, we conclude that the size $`L=40`$ can be regarded as sufficiently large to study the statistical properties of macroscopic quantities in our model. It is important to note here that the values of $`D`$ in the case $`E=10`$ are different from those in the case $`E=0`$.
Next, we define the conductivity $`\sigma `$ by adding a sufficiently small perturbativeg driving force $`ฯต`$ in the $`y`$-direction. This is realized by adding the term $`ฯตj_\mathrm{t}(๐ผ๐ผ^{})`$ to $`Q(๐ผ๐ผ^{})`$ in (1), where
$`j_\mathrm{t}(๐ผ๐ผ^{})`$ $``$ $`{\displaystyle \underset{i}{}}[\eta _i(1\eta _i^{})\eta _{i+(0,1)}^{}(1\eta _{i+(0,1)})`$ (5)
$``$ $`\eta _i^{}(1\eta _i)\eta _{i+(0,1)}(1\eta _{i+(0,1)}^{})].`$
Note that in the $`x`$-direction, the particles are still driven by $`E`$. Then, the averaged current $`\overline{J}_ฯต`$ in the $`y`$-direction is defined as
$$\overline{J}_ฯต\frac{1}{L}j_\mathrm{t}(๐ผ๐ผ^{})_\mathrm{s}^{E,ฯต}.$$
(6)
Using this $`\overline{J}_ฯต`$, the conductivity $`\sigma `$ is written
$$\sigma \underset{ฯต0}{lim}\frac{\overline{J}_ฯต}{ฯต}.$$
(7)
In the inset of Fig. 2, $`\sigma `$ is plotted as a function of the system size, $`L`$, in the cases $`E=0`$ and $`E=10`$ with $`\rho =0.5`$. Note that the values of $`\sigma `$ in the case $`E=0`$ are smaller than those in the case $`E=10`$, and that qualitatively, this difference is not the same as that seen for $`D`$.
We previously measured the intensity of density fluctuations $`\chi L\mathrm{}(\rho _{\mathrm{}}^2_\mathrm{s}^E(\rho _{\mathrm{}}_\mathrm{s}^E)^2)`$, where $`\rho _{\mathrm{}}_{i\mathrm{\Omega }_{\mathrm{}}}\eta _i/|\mathrm{\Omega }_{\mathrm{}}|`$ and $`\mathrm{\Omega }_{\mathrm{}}=\{(i_x,i_y)|1i_xL,L/2\mathrm{}/21i_yL/2+\mathrm{}/2\}`$. (See Fig. 4 in Ref. HSI .) Note that $`\mathrm{}`$ is chosen so that it satisfies $`\xi \mathrm{}L`$ where $`\xi `$ is a correlation length. Using these values of $`\chi `$, in Fig. 2, we plot $`\sigma T`$ as a function of $`D\chi `$ in the cases $`(\rho ,E)=(0.5,10),(0.4,10),(0.3,10),(0.5,3)`$ and $`(0.5,0)`$. Noting that the thin-dotted line represents $`D\chi =\sigma T`$, we find that even though the values of $`D`$, $`\sigma `$ and $`\chi `$ are different from those in the equilibrium case, along the direction transverse to the driving force $`E`$, the Einstein relation (2) is valid, within the precision of the numerical computations.
## Fluctuation-response relation:
We next study in the fluctuation-response relation, which is also a representative universal relation in the linear response theory. Again in this case, we focus on the properties of the system in the direction transverse to the driving force $`E`$.
First, we employ the same procedure as in the measurement of $`D`$ to introduce a time dependent response function $`R(t)`$. That is, we prepare the steady state under the perturbation $`V(i_y)`$ and then remove this perturbation at $`t=0`$. Because the profile of the coarse grained density $`\rho (i_y)`$ is changed by the removal of $`V(i_y)`$, we make this change explicit by defining $`R(t)`$ in the following form:
$$R(t)\frac{\widehat{\rho }(t)_E^V}{L\mathrm{\Delta }}.$$
(8)
We remark that the decaying behavior of $`R(t)`$ in the case $`E=10`$ is plotted as that of $`\widehat{\rho }(t)_E^V/\mathrm{\Delta }`$ in Fig. 1.
Next, we introduce the time correlation function of density fluctuations in the direction transverse to the driving force:
$$C(t)\widehat{\rho }(t)\widehat{\rho }(0)_\mathrm{s}^E.$$
(9)
In the inset of Fig. 3, $`\mathrm{ln}C(t)`$ is plotted as a function of time in the case $`E=10`$, with $`\rho =0.5`$ and $`L=32`$. For the NESS, $`C(t)`$, like $`R(t)`$ decays exponentially in time.
Here, in the equilibrium case ($`E=0`$), using $`R(t)`$ and $`C(t)`$, the fluctuation-response relation is given by
$$C(t)=TR(t).$$
(10)
In the NESS far from equilibrium studied here, because $`R(t)`$ and $`C(t)`$ exhibit a similar behavior, and because the fluctuation relation, which is essentially the same as $`C(0)=TR(0)`$, has previously been found to hold HSI , we conjecture that (10) is valid.
To demonstrate its validity explicitly, in Fig. 3, in the case $`E=10`$ with $`\rho =0.5`$ and $`L=32`$, $`R(t)`$ is plotted as a function of $`C(t)`$ over the interval $`0t800`$ MCS. It is seen that the slope is equal to $`1/T`$, within the precision of the numerical computations.
## Green-Kubo relation:
Finally, again considering the properties along the direction transverse to the external driving force (the $`y`$ direction), we investigate the validity of the Green-Kubo relation for NESS far from equilibrium.
Using the spatially-averaged current in the direction transverse to the external driving force, $`j_\mathrm{t}`$, defined in (5), we begin by the $`\tau `$ dependent current $`J^\tau `$, which represents the net number of particles that move in the $`y`$ direction during a time of $`\tau `$ MCS,
$$J^\tau \frac{1}{\tau L^2}\underset{k=1}{\overset{\tau L^2}{}}j_\mathrm{t}(๐ผ(k1)๐ผ(k)).$$
(11)
Then, using this expression for $`J^\tau `$, the intensity of the current fluctuations is defined by
$$B^\tau \frac{\tau L^2}{2}(J^\tau )^2_\mathrm{s}^E.$$
(12)
In Fig 4, $`B^\tau \tau `$ is plotted as a function of $`\tau `$ in the cases $`E=0`$ and $`E=10`$, with $`\rho =0.5`$ and $`L=32`$, respectively. It is seen that in the case $`E=10`$, the line fitted for small times (but much larger than the relaxation time of the current correlations) deviates slightly for large times, while in the case $`E=0`$, $`B^\tau \tau `$ and $`0.050\tau `$ are equal within the numerical precision for all times. This bending behavior of $`B^\tau \tau `$ might reflect the effect of a long time tail in this NESS.
In the case $`E=0`$, defining $`B`$ as the slope of $`B^\tau \tau `$, the Green-Kubo relation kubo can be written
$$B=\sigma T.$$
(13)
In the case $`E0`$, we define $`B`$ as the slope of $`B^\tau \tau `$ obtained from the fitting in the early time regime. With this definition, the size dependence of $`B`$ in the cases $`E=0`$ and $`E=10`$ with $`\rho =0.5`$ is plotted in the inset of Fig. 5. The difference between the values of $`B`$ in the cases $`E=0`$ and $`E=10`$ is qualitatively the same as that for $`\sigma `$.
Considering this similarity between $`\sigma `$ and $`B`$, in Fig. 5, we plot $`\sigma T`$ as a function of $`B`$ in the cases $`(\rho ,E)=(0.5,0),(0.5,10),(0.5,3),(0.4,10),(0.3,10)`$ with $`L=32`$. The Green-Kubo relation (13) is valid for the NESS considered here. However, the deviation seen in Fig. 5 is somewhat larger than that in Fig. 2.
## Summary:
In this paper, we have reported the results of numerical experiments on the two-dimensional DLG focusing on the properties along the direction transverse to the external driving force $`E`$. We find that the Einstein relation, the fluctuation-response relation and the Green-Kubo relation hold in the NESS far from equilibrium. (Note that the validity of the fluctuation relation for such a state was previously demonstrated in Ref. HSI .)
Compared with this validity of relations in the direction transverse to the external driving force, we remark the properties along the direction parallel to the external driving force. We studied a one-dimensional DLG, and found that the phenomena observed along the direction parallel to the external driving force seemed to be more complicated than those observed along the direction transverse to the external driving force HSIV .
We end with some discussion of the detailed balance of fluctuations, which has a deep connection with the validity of universal relations in the linear response regime near equilibrium states. In our DLG, the detailed balance condition for $`c(i,j;๐ผ)`$ does not hold in the case $`E0`$. However, the numerical confirmation of the universal relations presented here suggests the detailed balance of macroscopic fluctuations. We point out that, with regard to this topic, D. Gabielli et al. studied a stochastic model for which the detailed balance condition does not hold, and derived the Onsagerโs reciprocity, which is also the linear response relations for macroscopic quantities gabli .
The author acknowledges S. Sasa and H. Tasaki for discussions on NESS. This work was supported by a JSPS Research Fellowship for Young Scientists (Grant No. 1711222).
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# Physics Related with Co-moving Coordinate System
## I Introduction
The luminosity-distance v.s. red-shift relation of type Ia super-novae is the most direct evidence of dark energyโs existence see Riess98 ; Perlmutter98 ; Knop03 ; Tonry03 ; Riess04 for experimental literatures and Quintessence2 ; Phantom ; Phantom2 ; backReaction for theoretical explanations. The matter distribution power spectrum observed by SDSS Tegmark04 and cosmic microwave background anisotropy observed by WMAP WMAP03 is indirect evidence of dark energyโs existence. We will provide a new explanation for super-novaeโs luminosity-distance v.s. red-shift relation without assuming that the universe is accelerate-ly expanding. Our explanation will avoid almost all the main problems of standard cosmology before inflation is introduced into physics.
Our explanation resorts to new interpretation of physics related with co-moving coordinate system. As the first step, We would like to ask, if one focuses on a given direction of a non-perturbed universe, will he see an infinitely long, uniform and expanding galaxy line? If no, why? We consider only flat universe.
If yes, suppose this man/woman were put on galaxy $`O`$ and were asked to measure the recession velocity of galaxy $`B`$ and $`C`$, see FIG.1, what result will he/she get? $`(v,2v)`$ or $`(v,\frac{2v}{1+vv})`$? $`v`$ is the relative recession velocity between two nearest galaxies. We insist the second answer, i.e., we insist that (i) cosmological principle is a local statement; (ii) the definition of simultaneity can only be relativistic.
If the one dimensional system in Figure.1 is uniformly expanding, the metric is
$`ds^2=dt^2+a^2(t)dx_{co}^2,`$ (1)
when generalizing into (1+3)D space-time, we have
$`ds^2=dt^2+a^2(t)(dr_{co}^2+r_{co}^2[d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2])`$ (2)
Some standard cosmologists claim that, the generalization of $`(1+1)D(1+3)D`$ is ir-rationale. Because (1+1)D gravitation theory is topological, its $`G_{\mu \nu }==0`$, so no dynamical equations can be used to determine $`a(t)`$. However if we know the (1+1)D $`a(t)`$ somehow, the generalization to (1+3)D is rationale, because if we are considering a non-perturbed universe and if we are focusing on only a given direction, we will see the galaxy line in FIG. 1. Since the projection $`(1+3)D(1+1)D`$ involves only kinematic, it involves no dynamics.
So our claiming is: if we know the (1+1)D $`a(t)`$ somehow, generalization of eq(1) into eq(2) is rationale, because the generalization only involves kinematic, it involves no dynamics. Although the (1+1)D gravitation is topological, $`G_{\mu \nu }==0`$, the (1+3)D metric obtained by generalizing a pre-given (1+1)D metric has non-zero $`G_{\mu \nu }`$. So non-trivial dynamics can appear in the generalized (1+3)D space-time.
For example, if the system illustrated in FIG. 1 is expanding with zero acceleration, we can derive its metric โ(1+1)D formโ by pure kinematic method then generalize the results into (1+3)D space-time thus obtain the metric of an isotropic and homogeneous universe whose expansion has zero acceleration. The (1+1)D metric we obtained by pure kinematic method has identically zero $`G_{\mu \nu }`$, but the generalized (1+3)D metric has non-zero $`G_{\mu \nu }`$. So in the (1+3)D space-time, non-trivial dynamic appears. On the other hand, if we know the energy momentum tensor corresponding with an expanding universe with zero accelerations as priors, directly solving the (1+3)D dynamic equation will give us the same metric as that obtained by kinematic method of $`(1+1)D(1+3)D`$.
## II (1+1)D Expanding Universe with Zero Accelerations
Take the galaxy line in FIG.1 as our experimental labs. Suppose the system is expanding with zero-accelerations. Considering the following series
$`v_B=v;`$
$`v_C={\displaystyle \frac{v+v}{1+v^2}};`$
$`v_D={\displaystyle \frac{v+v_C}{1+vv_C}};`$
โฆ โฆ
$`v_X={\displaystyle \frac{v+v_{X1}}{1+vv_{X1}}};`$ (3)
$`|AB|=2a`$
$`|OC|=2a\sqrt{1v_B^2}`$
$`|BD|=2a\sqrt{1v_C^2}`$
โฆ โฆ
$`|X^{}X^+|=2a\sqrt{1v_X^2},`$ (4)
$`v`$, $`a`$ are locally measured relative recession velocity and distance between two nearest galaxies respectively. Since we consider only the non-accelerate-ly expanding universe, so $`a=vt`$. From series eqs(3) and (4) we get
$`v_X`$ $`={\displaystyle \frac{(1+v)^X(1v)^X}{(1+v)^X+(1v)^X}};`$ (5)
$`|OX|`$ $`a{\displaystyle \underset{N=0}{\overset{X}{}}}\sqrt{1v_N^2}`$ (6)
$`=l{\displaystyle _0^X}๐x\sqrt{1v_x^2}`$
$`={\displaystyle \frac{4a}{\text{ln}\frac{1+v}{1v}}}\left[\text{arctg}[({\displaystyle \frac{1+v}{1v}})^{\frac{X}{2}}]{\displaystyle \frac{\pi }{4}}\right],`$
From eq(6) using light velocity invariance principle, we can write down the metric of our (1+1)D non-accelerate-ly expanding universe as
$`ds^2=dt^2+{\displaystyle \frac{4v^2t^2}{(e^{\sigma x}+e^{\sigma x})^2}}dx^2`$ (7)
$`\text{where }\sigma ={\displaystyle \frac{1}{2}}\text{ln}{\displaystyle \frac{1+v}{1v}}`$ (8)
or
$`ds^2=dt^2+a^2(t)dx_{co}^2,a(t)=vt`$ (9)
$`\text{where}x_{co}={\displaystyle \frac{2}{\sigma }}(\text{arctg}[e^{\sigma x}]{\displaystyle \frac{\pi }{4}}).`$ (10)
We call the coordinate $`x`$ in eq(7) natural coordinate, while the coordinate $`x_{co}`$ in eq(9) co-moving coordinate. Natural coordinate ranges in $`(\mathrm{},\mathrm{})`$, but co-moving coordinate only ranges in $`(\frac{\pi }{2\sigma },\frac{\pi }{2\sigma })`$. If $`v0`$, natural coordinate coincides with co-moving coordinate. The co-moving coordinate definition of Standard cosmology emphasizes only one point: co-moving coordinate is a coordinate fixed on galaxies, the co-moving coordinate of a given galaxy does not vary as background universe expands. By this definition, natural coordinate is also co-moving coordinate.
But the difference between co-moving coordinate and natural co-ordinate is very important. Natural coordinate is related with physical coordinate through
$`x_{ph}={\displaystyle \frac{2vt}{\sigma }}(\text{arctg}[e^{\sigma x}]{\displaystyle \frac{\pi }{4}}).`$ (11)
While co-moving coordinate is related with physical coordinate through
$`x_{ph}=a(t)x_{co},a(t)=vt.`$ (12)
Although the co-moving coordinate defined in eq(12) is very similar to that of standard cosmology, it has completely different interpretation from that of standard cosmology.
By standard cosmologyโs definition $`x_{ph}=a(t)x_{co}`$, if we have a photon emitted at $`(t,x_{co})`$ and detected at $`(t_0,0)`$, the red-shift of this photon is $`(1+z)=\frac{a(t_0)}{a(t)}`$, it has nothing to do with the co-moving coordinate of source galaxy. But in the co-moving coordinate definition of eq(12), for the same photon, the red-shift is
$`(1+z)=\sqrt{{\displaystyle \frac{1+v_x}{1v_x}}}=e^{\sigma x}`$ (13)
It is completely determined by the co-moving coordinate of source galaxy but has nothing to do with scale factor!
If we accept eqs(11)+(13), then even without assuming that the universe is accelerate-ly expanding, we can give the observed luminosity-distance v.s. red-shift relation of super-novaes a very beautiful explanation. Of course, if one would like to, he can replace the velocity $`v`$ in eqs(7), (9), (11) and (12) with a time dependent function $`v(t)=v+pt+\frac{1}{2}qt^2+\mathrm{}`$, and the factor $`a(t)=vt+\frac{1}{2}pt^2+\frac{1}{3!}qt^3+\mathrm{}`$, thus obtain appropriate relations in a accelerate-ly expanding universe. In this case eq(13) will be changed so that more free parameters enter, hence more precise fitting with experiments can be obtained.
## III Generalization into (1+3)D Space-time
Generalizing eqs(7) or (9) into (1+3)D space-time, we obtain
$`r_{ph}`$ $`={\displaystyle \frac{2vt}{\sigma }}(\text{arctg}e^{\sigma r}{\displaystyle \frac{\pi }{4}})`$ (14)
$`ds^2`$ $`=dt^2+{\displaystyle \frac{v^2t^2}{\mathrm{cosh}^2\sigma r}}(dr^2+r^2[d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2])`$ (15)
$`r_{ph}`$ $`=a(t)r_{co},a(t)=vt,0r_{co}{\displaystyle \frac{\pi }{2\sigma }}`$ (16)
$`ds^2`$ $`=dt^2+a^2(t)(dr_{co}^2+r_{co}^2[d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2])`$ (17)
The parameter $`v`$ now should be understood as the average recession velocity between two nearest galaxies. If we take a limit $`v0`$, co-moving coordinate reduce to natural coordinate, eq(17) becomes (15). Note, we consider only galaxies which are performing Hubble recessions relative to each other. We do not consider galaxies bounded in the galaxy clusters.
Just the same as (1+1)D case, eq(17) is very similar to standard cosmologyโs FRW metric, but the two has completely different physical interpretations. E.g. if we have a photon emitted at $`(t,r_{co},\theta ,\varphi )`$ and detected at $`(t_0,0,\theta ,\varphi )`$. By standard cosmology, the red-shift of this photon is $`(1+z)=\frac{a(t_0)}{a(t)}`$; but by our explanation of eq(17), the red-shift is
$`(1+z)=\sqrt{{\displaystyle \frac{1+v_r}{1v_r}}}=e^{\sigma r}`$ (18)
Although our starting point, the (1+1)D metric eqs(7) and (9) are topological theory, it has no dynamics. When we generalize them into (1+3)D case, eqs(15) or (17), non-trivial dynamics appears. By Einstein equation, we can calculate the energy momentum tensor corresponding with them. The results are respectively
$`8\pi GT_{\mu \nu ,na}`$ $`=G_{\mu \nu ,na}=\text{diag}`$ (19)
$`\{{\displaystyle \frac{(6v^2r+5\sigma ^2r\sigma ^2r\mathrm{cosh}[2\sigma r]+4\sigma \mathrm{sinh}[2\sigma r])}{2v^2rt^2}}`$
$`,(\sigma ^2+v^2)\text{sech}^2[\sigma r]+\sigma (\sigma +{\displaystyle \frac{2\mathrm{tanh}[\sigma r]}{r}})`$
$`,r^2(\sigma ^2+v^2)\text{sech}^2[\sigma r]+\sigma r\mathrm{tanh}[\sigma r],`$
$`\mathrm{sin}^2\theta [r^2(\sigma ^2+v^2)\text{sech}^2[\sigma r]+\sigma r\mathrm{tanh}[\sigma r]]\}`$
$`8\pi GT_{\mu \nu ,co}`$ $`=G_{\mu \nu ,co}`$ (20)
$`=\text{diag}\{{\displaystyle \frac{3}{t^2}},v^2,v^2r_{co}^2,v^2r_{co}^2\mathrm{sin}^2\theta \}`$
If we know the energy momentum tensor describing the cosmological fluid is $`T_\nu ^\mu =\text{diag}\{\rho ,\frac{1}{3}\rho ,\frac{1}{3}\rho ,\frac{1}{3}\rho \}`$ in the co-moving coordinate, i.e. $`T_{\mu \nu }`$ expressed in eq(20) times $`g^{\mu \nu }`$, then starting from a general ansaltz $`ds^2=dt^2+a(t)(dr_{co}^2+r_{co}^2d\mathrm{\Omega }^2)`$, using Einstein equation, we can also derive out the function form of $`a(t)=vt`$. This is just the routine of standard cosmology. But this routine slides over all physics related with the definition of co-moving coordinate system, see eq(18) and the related remarks.
From eq(20), we can see that in the co-moving coordinate system, for a non-accelerately expanding universe, its cosmological fluid has pressure $`p=\frac{1}{3}\rho `$. Note it is $`T_\nu ^\mu `$, not $`T_{\mu \nu }`$, that is directly related with energy density and pressure, $`T_\nu ^\mu =\rho u^\mu u_\nu +p(u^\mu u_\nu +\delta _\nu ^\mu )`$. We have two reasons to accept this negative pressure.
The first reason is, we can think it originates from dark energy, e.g., $`T_\nu ^\mu `$ = $`(\rho ,\frac{1}{3}\rho ,\frac{1}{3}\rho ,\frac{1}{3}\rho )`$ = $`(\frac{2}{3}\rho ,0,0,0)`$ \+ $`(\frac{1}{3}\rho ,\frac{1}{3}\rho ,\frac{1}{3}\rho ,\frac{1}{3}\rho )`$. Note, in standard cosmology, to prevent the expansion of universe from decelerating, dark energy must be included in the total energy menu of the universe.
The second reason is, even in a universe containing only matters, negative pressure can appear as a result of gravitations. Imagine an infinitely long uniform galaxy line, if one of the composite galaxies is less weighted than others, then all galaxies on the left hand side of this less weighted galaxy will collapse and move to the left, while all galaxies on the right hand side of this less weighted galaxy will collapse and move to the right. So, the less weighted galaxy receives gravitations which have intentions to split it into two parts. This intention can be understood as the origin of negative pressure. A less weighted galaxy is just an auxiliary object to illustrate the effects negative pressure. When all galaxies are equal weighted, negative pressure also exists as a result of gravitations.
If we accept the first reason, i.e., negative pressure originates from dark energy, then we will have to accept that our metric eqs(15) and (17) only describe our universe in a very short period of time. Because $`\rho _ma^3(t)`$, while $`\rho _{de}a^{3(1+w)}(t)`$. If $`\rho _m\rho _{de}`$ today, then in the far past, $`\rho _{de}`$ must be much less than $`\rho _m`$, so negative pressure provided by dark energy will not be able to prevent the universe from decelerate-ly expanding. If we accept the second reason, i.e. dark energy originates from gravitations between different parts of the universe, then our metric eqs(15) and (17) can be used to describe the universe in any eras when the gravitation is the main interaction between different parts of the universe.
Of course, at very early times, galaxies do not exist, so the parameter $`v`$ cannot be understood as the average recession velocity between two nearest galaxies, but according to the hierarchical clustering scenario, before galaxies appear, stars exist, before stars appear, nucleon exists, before nucleon appears, electron and protons exists, โฆ . So, as long as we accept that inter-gravitations among different parts of the universe produce negative pressures, while the average velocity of relative recession between two nearest composite object is $`v`$, then eqs(15) and (17) can be used to describe our universe at times as early as big-bang nucleon synthesis, even primordial singular point. While God, need only calculate and assign value to one parameter $`v`$, so that when the universe is 137Gyr old, human appears.
## IV Observations of Super-novae
From theoretical aspects, the basis of eqs(15) and (17) is very simple and concrete, (i) cosmological principle is a local statement; (ii) the definition of simultaneity can only be relativistic. However, will experimental observations support it? The luminosity-distance v.s. red-shift relation of super-novaes is the most direct evidence that the universe is accelerate-ly expanding. But this statement is based on ignoring of physics related with co-moving coordinate system discovered in this paper. We will show that when considering physics related with co-moving coordinate system, the observational result can be explained even without assuming that our universe is accelerate-ly expanding.
If we consider physics related with co-moving coordinate system, the red-shift of photons coming from distant galaxies will be changed remarkably comparing with standard cosmology. When the universe is assumed expanding with zero acceleration, photons emitted from a super-novae at position $`(t,r,\theta ,\varphi )`$ have red-shift
$`(1+z)=\sqrt{{\displaystyle \frac{1+v_r}{1v_r}}}=e^{\sigma r}.`$ (21)
Considering Lorentz dilating, the photons emitted in period $`\delta t_1`$ can only reach us in period $`\delta t_1e^{\sigma r}`$. So we get the luminosity-distance v.s. red-shift relation as
$`d_l=(1+z){\displaystyle \frac{2vH_0^1}{\sigma }}[\text{arctg}(1+z){\displaystyle \frac{\pi }{4}}]`$ (22)
Please refer to SWeinberg , section 14.4, eqs(14.4.11-14) for detailed derivation of eq(22).
From FIG. 2, we can see that when considering physics related with co-moving coordinate system, even without assuming that our universe is accelerate-ly expanding, theoretical predictions are very close to predictions of $`\mathrm{\Lambda }`$CDM cosmology. From best fitting observational results of Riess04 , we get $`v=0.79/3000`$, $`H_0=60\text{km/(s}\text{Mpc)}`$, $`\chi ^2=303`$ (186Golden+Silver sample) or $`v=0.899/3000`$, $`H_0=`$$`60\text{km/(s}\text{Mpc)}`$, $`\chi ^2=237`$ (157Golden sample).
Only judging from numerical fitting qualities, our prediction eq(22) may be not as good as standard cosmology. But our theoretical frame-work has only two free parameters, $`v`$ and $`H_0`$, while standard cosmology actually uses three parameters, $`\mathrm{\Omega }_{m0}`$, $`H_0`$ and $`w`$. The superiority of our frame-work over standard cosmology is mainly on theoretical aspects.
Standard cosmology does not give any explanation of negative pressureโs producing mechanism, so the equation of state coefficient $`w`$ of dark energy must be counted as a free parameter. But we provide a possible negative pressure producing mechanism, it can produce $`p=\frac{1}{3}\rho `$. We will discuss this problem in more details in the discussion section. Since we do not resort to dark energies to explain the luminosity-distance v.s. red-shift relation of super-novaes, our universe contains only matters today, so our cosmological frame-work has no coincidence problem Quintessence1 .
Since in our cosmological frame-work, universe expands with zero acceleration, the size of the observable universe is always equal to the particle horizon of the universe. So our frame-work has no horizon problem. Since our cosmological frame-work has no horizon problem, quantum fluctuations inside the horizon will provide the primordial seeds for latter structure formations. So our frame-work has no primordial structure formation seeds problem Guth81 ; Linde82 ; EkpyroticUniverse .
Considering physics related with co-moving coordinate system, the global topology of the universe is not related with the energy density of the universe through a simple Friedmann equation $`\frac{\dot{a}^2}{a^2}+\frac{k}{a^2}=\frac{8\pi G}{3}\rho _{tot}`$. We have not found the metric of a closed/open universe by kinematical method. Probably, eq(17) is the only solution of the real universe. If that is the case, our cosmological frame-work has no flatness problem at all.
## V Discussions
There are two worries about our considering of physics related with co-moving coordinate system. The first is, since we use special relativity velocity addition rules to calculate recession velocity of galaxies on a given observational direction, some people worry that our theory will contradict the basic fact of Hubbleโs discovery, the recession velocity of a galaxy is proportional to the distance the galaxy being away from us, $`vH_0x_{ph}`$. First let me explain that, even for this basic fact, different standard cosmologists could give us different interpretations.
The first class standard cosmologists say that this is an empiric formulae only valid at low red-shift. Because if it is valid at very high red-shift, or on very large physical distances, super-light velocity of recessions would appear, which is anti-relativity. The second class standard cosmologists say that, $`vH_0x_{ph}`$ is a basic principle valid on any scales; super-light recession velocity on super-horizon scales introduces no problem, so is allowed. We support the first class of standard cosmologists, i.e., $`vH_0x_{ph}`$ is an empiric formulae only valid on low red-shift or small (compare with observational horizon of the universe) scales.
Still take the one-dimensional galaxy line in FIG. 1 as our experimental labs. Consider the recession velocity and the physical distance of a galaxy located at $`(t,x)`$ relative to us, $`x`$ is the natural coordinate of that galaxy,
$`v_x`$ $`={\displaystyle \frac{e^{\sigma x}e^{\sigma x}}{e^{\sigma x}+e^{\sigma x}}};`$ (23)
$`x_{ph}`$ $`={\displaystyle \frac{4a}{\text{ln}\frac{1+v}{1v}}}\left[\text{arctg}[({\displaystyle \frac{1+v}{1v}})^{\frac{X}{2}}]{\displaystyle \frac{\pi }{4}}\right],`$ (24)
Obviously, only when $`v_x<<1`$, i.e. $`\sigma x<<1`$, $`v_xx_{ph}`$. Obviously, regardless how large is $`x`$, the recession velocity $`v_x`$ cannot be larger than $`1`$, the light velocity.
The second worry about our cosmological picture is, since we assume that the average value of relative recession velocity between two nearest galaxies is time independent. The Hubble parameter is also time independent. This is just an illusion. Still take the one-dimensional galaxy line as our experimental labs. Obviously, since the distance between two nearest galaxies increase linearly with time, Hubble parameter, decreases as $`t^1`$ as time passes by. This is the same as standard cosmologyโs matter/radiation dominated erasโ Hubble parameter evolution rules. So if we use eq(17) to trace back the history of our universe, we will not get result inconsistent with Big Bang Nucleon-synthesis of standard cosmology.
Our final discussion is about the negative pressure $`p=\frac{1}{3}\rho `$โs producing mechanism. As the first step, let me ask if we have a one-dimensional infinitely long uniform galaxy line, and if the system is at rest initially, will it collapse at self-gravitations? Professor Ed. Witten once told me, Einstein contemplated similar questions. He considered a three-dimensional uniform lattice system, by poisson equation $`^2\varphi =\rho `$ (in general relativity, there are similar equations which will give us the same conclusions), the system has a solution $`\varphi =\frac{1}{2}\rho x^2`$. So for any galaxy not on the $`x=0`$ plane, it will receive a force pointing to that plane. As a result the system will collapse to that plane. Of course, the system could also has solution like $`\varphi =\frac{1}{2}\rho (xx_0)^2`$, which means that the system should collapse to the $`x=x_0`$ plane. So Einstein concludes that an isotropic and homogeneous matter dominated universe cannot have static solution. This is why Einstein introduced cosmological constant into his basic equations to get static solutions, as early as before Hubble discovered that our universe is expanding.
However, we wish to express a modest suspicion that, Einstein may ignore an important thing. In a infinitely long uniform galaxy line, inter-gravitations among different galaxies can produce negative pressures. Imagine that, there is a galaxy in the line containing less matters comparing with other galaxies. In this case, galaxies on the left hand side of this less weighted galaxy will collapse and move to the left, while galaxies on the right hand side of this less weighted galaxy will collapse and move to the right. So the less weighted galaxy will receive gravitations from both sides, which have intentions to split this galaxy into two parts. This intention can be understood as origins of negative pressure $`p=\frac{1}{3}\rho `$. Then why is the equation of state coefficient $`\frac{1}{3}`$?
Before answering this question, let us first imagine that, if what we illustrated in FIG. 1 is not a galaxy line, but an electron-line. Will the system expand at self-repulsion? According to the same poisson equation analysis of Einstein, the system should not expand at self-repulsion, but should collapse at self-repulsions! This is un-acceptable. So, maybe Einstein, and almost all standard cosmologists since Einstein, were cheated by their intuitions: an infinitely long uniform galaxy line will collapse at self-gravitations. They analyzed this problem by first ignoring pressures originated from the inter-gravitations(or static electronic repulsions) among the composite objects of the system then using the so-called dynamic equations (Einstein equation or Poisson equation) so get their conclusions.
Our point of view is, to answer the question that, will an infinitely long uniform galaxy line collapse at self-gravitations, or an electron line expanding at self-repulsions? We should not use dynamic equations at the condition of ignoring pressures originated from inter-actions among different composite objects. Otherwise, we will get un-acceptable conclusions, e.g., an infinitely long uniform electron-line will collapse at self-repulsions.
We think the reasonable conclusion should be, from symmetry analysis, any galaxy on the line receives gravitations from both sides. The two-side gravitations cancel each other, so any galaxy on the line will not run close to its neighbors, i.e., the system will not collapse or expanding at self-gravitations or self-repulsions. If we insist this analysis, then an initial expanding non-perturbed universe will keep expanding at the same speed for ever. For such a universe, we have used kinematic method and derived its metric in eq(17), by Einstein equation, the energy momentum tensor corresponding with metric has pressure $`p=\frac{1}{3}\rho `$.
We must claim that, when we say Einstein and his following standard cosmologists were cheated by their intuitions, we only want to express our modest suspicion. We have received many many criticisms and lampoons for our suspicion. But we think this suspicion is worth being kept in mind. After all, this suspicions has given us a possible explanation of the luminosity-distance v.s. red-shift relation of type Ia super-novae. We wish further exploration of this suspicion, for example, perturbing eq(15) or (17) and studying the structure formation or cosmic micro-wave background anisotropy problem then comparing with experimental observations such as SDSS Tegmark04 and WMAP WMAP03 will tell us whether our suspicion can be fact or not.
## VI Conclusions
We derive the metric of an expanding universe with zero accelerations by pure kinematic method. By doing so we expatiate physics related with co-moving coordinate system in details. The most important discovery of our study is, in an expanding universe with zero accelerations, the red-shift of photons from distance galaxies is determined by the co-moving coordinate of the source galaxy instead of the scale factor ration $`\frac{a(t_0)}{a(t)}`$. Our discovery is consistent with the current observed super-novaesโs luminosity-distance v.s. red-shift relations.
We also discuss that, an expanding universe with zero accelerations has no horizon problem, (probably)no flatness problem, no primordial structure formationโs seed problem. By Einstein equation, we find that to assure an expanding universe with zero accelerations, then energy momentum tensor of the underlying cosmological fluid must have $`p=\frac{1}{3}\rho `$. We discuss that such a negative pressure can originate from the inter-gravitations among different composite objects of the universe โ the galaxies.
Acknownedgement
Originally, this paper appears as an answering letter to criticisms on our works CosmoSDSF . When we finish the first version of that paper, we send it to professor E. Witten, G. โt Hooft, P. J. Steinhardt and other peoples for comments and criticisms. They read our paper and give comments seriously. We thank them very much for their comments or criticisms on our work in that paper. Their reactions encourage us very much.
The current version of this paper include results of discussions with professor L. Liu, S.-y Pei. They inquire me to give a talk on the topic discussed in CosmoSDSF at Beijing Normal University. When I finish the demonstrating document for the talk, I find the document itself may express my ideals more clearly than the original paper. So I decide to update the answering letter with the current paper โin factโ the demonstrating document of talk to be given in BNU.
Although I ask so many people to read my paper and give comments or criticisms on it, and they indeed do so. This does not mean that they agree with me on my opinions. So none of these people is to take response for the errors in the paper. But if there is any reasonable points in the paper, I must owe the credit to all of them.
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# R-charged ๐ดโข๐โข๐โ
black holes and large N unitary matrix models
## 1 Introduction
The AdS/CFT correspondence implies that the phases of string theory can be studied by studying those of the dual gauge theory(). In the case of type $`IIB`$ string theory in $`AdS^5\times S^5`$ the phases can be studied using the large N limit of a unitary matrix model()<sup>1</sup><sup>1</sup>1Phases of large N gauge theory is also discussed in . The unitary matrix is the finite temperature Polyakov loop which does not depend on the points of $`S^3`$. This fortunate circumstance is due to the fact that in the Hamiltonian formulation, N=4 SYM theory at a given time slice, is defined on the compact space $`S^3`$ and the $`SO(6)`$ scalars are massive because of their coupling to the curvature of $`S^3`$. These facts imply that, in principal, one can integrate out almost all the fields and obtain an effective theory of the zero mode of the gauge potential $`A_0`$.
Using this method a detailed correspondence of the critical points of the gauge theory effective lagrangian and the critical points of supergravity (discussed by Hawking and Page) can be constructed at the leading order of the $`1/N`$ expansion. These are $`AdS_5`$ and the small and big black holes (as we refer to these as SSB and BBH.) It turns out that in the gauge theory these critical points are in the gaped phase, where the density of eigenvalues vanishes in a finite arc of the circle around which the eigenvalues are distributed. The closing of the gap, corresponds to the Gross-Witten(GW) phase transition(). In a window around this transition, the supergravity description of string theory is likely to smoothly cross over into a description in terms of heavy string modes. <sup>2</sup><sup>2</sup>2In , besides the vicinity of the Gross-Witten transition, authors also studied $`1/N`$ corrections and presented formulas for the partition function in the vicinity of blackhole nucleation and the Hagedorn transition.
In this paper we extend the discussion of the correspondence between R-charged $`AdS_5`$ blackholes(), and the effective unitary matrix model(). R-charged black holes are known to have a rich phase structure in the canonical and grand canonical ensemble. In the canonical ensemble the fixed charge constraint, contributes an additional logarithmic term $`\mathrm{log}(TrUTrU^{})`$ involving the order parameter, to the gauge theory effective action. This term is crucial for matching with supergravity. We analyze the implications of this term in the large N limit and compare with the various supergravity properties like the existence of only blackhole solutions in the canonical ensemble and also the existence of a point of cusp-catastrophe in the phase diagram.
The plan of this paper is as follows. In section 2 we give a brief review of charged $`AdS_5`$ blackholes. In section 3 and section 4 we discuss the effective action of the gauge theory at zero and small coupling, in the fixed charged sector. At zero coupling there is exactly one saddle point and the value of $`TrUTrU^1`$ at the saddle point is always non-zero. For a small positive coupling there are two stable and one unstable saddle points, all with a non-zero value of $`TrUTrU^1`$. They merge at the $`GW`$ point. In section (5) we discuss the model effective action at strong coupling. Here too, there are three saddle points,two stable(I,III) and one unstable(II). In the region $`\rho >\frac{1}{2}`$, $`I`$ and $`III`$ can be identified with a stable small blackhole and stable big blackhole respectively. Saddle point $`II`$ is identified with the small unstable black hole. The merging of saddle points leads to critical phenomenon whose exponents can be calculated and shown to agree with supergravity. This is discussed in section 6. We have also calculated the $`o(1)`$ part of the partition function near the critical point.
## 2 R-charged blackholes in $`AdS_5`$ and critical phenomena
The R-charged $`AdS_5`$ black hole and relevant phase structure were discussed by A. Chamblin, R. Emparan, C. V. Johnson and R. C. Myers (,). Here we review their result. The EinsteinโMaxwellโantiโdeSitter (EMadS<sub>n+1</sub>) action may be written as
$$I=\frac{1}{16\pi G}_Md^{n+1}x\sqrt{g}\left[\stackrel{~}{R}F^2+\frac{n(n1)}{R^2}\right],$$
(1)
where $`\stackrel{~}{R}`$ is the Ricci scalar and $`R`$ is the characteristic length scale of $`AdS`$. The metric of the ReissnerโNordstrรถmโantiโdeSitter (RNadS) solution is given in static coordinates by
$$ds^2=(1\frac{m}{r^{n2}}+\frac{q^2}{r^{2n4}}+\frac{r^2}{R^2})dt^2+\frac{dr^2}{1\frac{m}{r^{n2}}+\frac{q^2}{r^{2n4}}+\frac{r^2}{R^2}}+r^2d\mathrm{\Omega }_{n1}^2,$$
(2)
The parameter $`q`$ is proportional to the charge
$$Q=\sqrt{2(n1)(n2)}\left(\frac{\omega _{n1}}{8\pi G}\right)q$$
(3)
and $`m`$ is proportional to the ADM mass $`M`$ of the blackhole.
$$M=\frac{(n1)\omega _{n1}}{16\pi G}m$$
(4)
$`\omega _{n1}`$ is the volume of the unit $`(n1)`$โsphere, and the gauge potential is given by
$$A_0=\left(\frac{1}{\sqrt{\frac{2(n2)}{n1}}}\frac{q}{r^{n2}}+\mathrm{\Phi }\right)$$
(5)
where $`\mathrm{\Phi }`$ is the electrostatic potential difference between the black hole horizon and infinity.
For n=4 the solution (2) can be considered as a rotating black hole in $`AdS_5\times S^5`$ with angular momentum in the internal space $`S^5`$.<sup>3</sup><sup>3</sup>3The EMadS system described here may be thought as the dimensional reduction of $`AdS_5\times S^5`$ to $`AdS_5`$. Generally introducing angular momentum in $`S^5`$ will distort $`S^5`$. This distortion is not taken into in the EMadS reduction . We thank S.Trivedi for pointing this out to us. However these details do not change our main consideration. The symmetry group of $`S^5`$ is $`SO(6)`$ and the black hole we are discussing has equal $`U(1)`$ charges for all the three commuting $`U(1)`$ subgroups of $`SO(6)`$, the R-symmetry group of the $`N=4`$ $`SYM`$. Hence we are dealing with a system which has the same chemical potential $`\mu `$ for all three $`U(1)`$ charges in the grand canonical ensemble or equivalently three fixed equal U(1) charges in the canonical ensemble.
### 2.1 Equation of state
In order to discuss the thermodynamics, we consider the Euclidean continuation ($`ti\tau `$) of the solution, and identify the imaginary time period $`\beta `$ with the inverse temperature. Using the formula for the period, $`\beta =4\pi /V^{}(r_+)`$ (for a review see ), we get
$$\beta =\frac{4\pi l^2r_+^{2n3}}{nr_+^{2n2}+(n2)l^2r_+^{2n4}(n2)q^2l^2}.$$
(6)
This may be rewritten in terms of the potential as:
$$\beta =\frac{4\pi l^2r_+}{(n2)l^2(1c^2\mathrm{\Phi }^2)+nr_+^2}.$$
(7)
The condition for euclidean regularity used to derive (6) is equivalent to the condition that the black hole is in thermodynamical equilibrium. The equation (6) may therefore be written as an equation of state $`T=T(\mathrm{\Phi },Q)`$. From this equation of state we see that for fixed $`\mathrm{\Phi }`$ we get two branches, one for each sign, when the discriminant under the square root is positive. For fixed $`Q`$, $`T(\mathrm{\Phi })`$ has three branches for $`Q<Q_{crit}`$ (Let us call them $`I`$, $`II`$, $`III`$) and one for $`Q>Q_{crit}`$. The critical charge is determined at the โpoint of inflectionโ by $`\left(Q/\mathrm{\Phi }\right)_T=\left(^2Q/\mathrm{\Phi }^2\right)_T=0.`$
The qualitative features of $`\beta (r_+)`$ for varying $`q`$ are shown in Fig 1.
There is a critical charge, $`q_{\mathrm{crit}}`$, below which there are three solutions for $`r_+`$ for a range of values of $`T`$, corresponding to small ($`I`$), unstable($`II`$), and large ($`III`$) black holes. For fixed $`q<q_{crit}`$ only branch $`I`$ is available at low temperatures. At $`T=T_{02}(q)`$, there is a nucleation of two new solutions, $`II`$ the unstable small black hole solution, and $`III`$ the stable big black hole solution. As the temperature is increased further, the black hole $`II`$ approaches black hole $`I`$ and at $`T=T_{02}(q)`$ the two solutions merge.
As $`q`$ is increased further, $`T_{02}`$ increases, whereas $`T_{01}`$ decreases. At $`q=q_{crit}`$ we have $`T_{crit}=T_{01}=T_{02}`$. At $`q_{crit}`$ and $`T_{crit}`$ all three solutions merge. In the language of catastrophe theory this is a cusp catastrophe. As we increase $`q`$ beyond $`q_{crit}`$ there will be just one solution for all temperatures $`T`$.
We wish to take note of some properties of the phase diagram.
1. Thermal $`AdS`$ is not a solution and all three branches of the solution represent black holes.
2. There exists a critical point where the three solutions of the system merge. It is a point of fold catastrophe.
### 2.2 Critical Phenomena
The critical point ($`Q_{cit}`$,$`T_{crit}`$) may be approached from various directions in the parameter space. If we set $`T=T_{crit}`$, then the equation determining $`(rr_{crit})`$ takes the form ($`r_{crit}`$ is the radius of the critical black hole)
$`(rr_{crit})^3=๐(QQ_{crit})`$ (8)
$`๐`$ is a numerical constant. The critical exponent here is $`\frac{1}{3}`$, since
$$(rr_{crit})(QQ_{crit})^{\frac{1}{3}}$$
(9)
As discussed in (Fig16), the critical point may also be approached through the coexistence line in the parameter space. The coexistence line is the line with the property
$$S_I=S_{III}$$
for the parametric range $`q<q_{crit}`$. It is the line where the Hawking-Page (first order) transition from the small black hole to the big black hole takes place. As we approach the critical point through this line, we have the relation
$`(r_Ir_{II})(TT_{crit})^{\frac{1}{2}}`$ (10)
In the following we will present an understanding of these properties. Before we do the matching with supergravity we would like to present a discussion of the gauge theory in the limit when $`\lambda =g^2N=0`$ and also when $`\lambda <<1`$. In these cases we of course can not compare with supergravity which requires $`\lambda >>1`$.
## 3 Free YM theory
### 3.1 Effective action with chemical potential
In this section we briefly review (see ,) the effective action for a free $`SU(N)`$ Yang-Mills theory (with adjoint matter) on a compact manifold $`\mathrm{}`$ in the large $`N`$ limit. The basic idea is to integrate out all fields in the theory except for the zero mode of the Polyakov line. The partition function is then reduced to a single unitary matrix integral.
Expanding all fields in the gauge theory in terms of harmonics on $`\mathrm{}`$, the theory reduces to a zero dimensional problem of free $`N\times N`$ Hermitian matrices
$`={\displaystyle \frac{1}{2}}{\displaystyle \underset{a}{}}Tr\left[(D_tM_n)^2\omega _n^2M_n^2\right]`$ (11)
The sum in (11) )is over all field types and their Kaluza-Klein modes on $`\mathrm{}`$. $`\omega _n`$ is the frequency of each mode. The covariant derivative in (11) is
$$D_tM_n=_tM_ni[A_0,M_n]$$
(12)
$`A_0`$ comes from the zero mode (i.e. the mode independent of coordinates on $`\mathrm{}`$) of the time component of the gauge field and is not dynamical. The partition function of the theory at finite temperature can be written as a unitary matrix integral by integrating out all fields in (11) except for $`A_0`$. Hence we have
$`Z`$ $`={\displaystyle DU\underset{n}{}\left(\mathrm{det}_{adj}\left(1ฯต_ne^{\beta \omega _n}U\right)\right)^{ฯต_n}}`$ (13)
where $`U=\mathrm{exp}(i\beta A_0)`$ is a $`U(N)`$. $`\mathrm{det}_{adj}`$ denotes the determinant in the adjoint representation and $`ฯต_n=1`$ ($`1`$) for bosonic (fermionic) $`M_n`$. The above equation can be expressed as
$`Z={\displaystyle }DU\mathrm{exp}({\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{z_n(\beta )}{n}})TrU^nTrU^n)`$ (14)
where
$`z_n(\beta )=z_B(n\beta )+(1)^{n+1}z_F(n\beta ).`$ (15)
Here $`z_B(\beta ),z_F(\beta )`$ are the single particle partition functions of the bosonic and fermionic sectors respectively ( see for the explicit formulas for $`z(\beta )`$ in various gauge theories).
If we introduce a chemical potential $`\mu =(\mathrm{log}m)`$, the formula for the partition function changes to
$$Z[\beta ,\mu ]=DUe^{(_{n=1}^{\mathrm{}}\frac{z_n(\beta ,\mu )}{n}TrU^nTrU^n)}$$
(16)
where
$$z_n(\beta ,\mu )=z_B(n\beta ,n\mu )+(1)^{n+1}z_F(n\beta ,n\mu )$$
(17)
$`z_B(\beta ,\mu ),z_F(\beta ,\mu )`$ are now the single particle partition functions with chemical potential $`\mu `$. Hence
$`z_B(\beta ,\mu )`$ $`=`$ $`{\displaystyle \underset{bosons}{}}\mathrm{exp}(E_i\beta +Q_i\mu )`$ (18)
$`z_F(\beta ,\mu )`$ $`=`$ $`{\displaystyle \underset{fermions}{}}\mathrm{exp}(E_i\beta +Q_i\mu )`$ (19)
$`Q_i`$ is the charge of the state whose energy is $`E_i`$ (see ).
As we have already mentioned, we wish to describe a system which has the same chemical potential $`\mu `$ for all three $`U(1)`$ charges of the R-symmetry group $`SO(6)`$ in grand canonical ensemble. Equivalently, we can work with fixed and equal values of the $`U(1)`$ charges in the canonical ensemble.
Let us for simplicity confine ourselves to the bosonic sector of the $`N=4`$ SYM theory. The gauge fields have no R-charge. The six scalars $`\varphi _i`$ $`(i=1,\mathrm{},6)`$ can be grouped in pairs of two. We define
$`\varphi ^+=\varphi ^1+i\varphi ^2,`$ $`\varphi ^{}=\varphi ^1i\varphi ^2`$ (20)
(We can similarly define complex fields for the other two pairs.) $`\varphi ^\pm `$ have charge $`\pm 1`$ for each of the three commuting $`U(1)s`$ of $`SO(6)`$. Hence, if we consider the single particle partition function for these fields, it will be
$$z[x,\mu ]=(exp(+\mu )+exp(\mu ))z[x,0]/2=\mathrm{cosh}(\mu )z[x,0]$$
(21)
where $`z[x,0]`$ is the single particle partition function without any chemical potential.
### 3.2 Canonical Ensemble
We will now discuss the free gauge theory partition function for a canonical ensemble with constant charge $`Q_0`$, by introducing a delta function $`\delta (\widehat{Q}Q_0)`$ in the functional integration of the gauge theory. $`\widehat{Q}=Q[\varphi ]`$ is the corresponding functional for the charge which we want to keep fixed. In our case $`\widehat{Q}`$ is just the functional for R-charge in gauge theory.
The fixed charge partition function is defined by
$`Z(\beta ,Q_0)`$ $`=`$ $`{\displaystyle DXe^{_0^\beta S[X]}\delta (\widehat{Q}Q_0)}`$ (22)
$`=`$ $`{\displaystyle DXe^{_0^\beta S[X]}\mathrm{exp}(i\mu \widehat{Q})\mathrm{exp}(i\mu Q_0)๐\mu }`$ (23)
$`=`$ $`{\displaystyle ๐\mu \mathrm{exp}(i\mu Q_0)(DXe^{_0^\beta S[X]}e^{i\mu \widehat{Q}})}`$ (24)
$`=`$ $`{\displaystyle ๐\mu \mathrm{exp}(i\mu Q_0)DU\mathrm{exp}(\mathrm{\Sigma }z_n[\beta ,i\mu ]TrU^nTrU^n)}`$ (25)
where $`z_n[\beta ,i\mu ]=z_n^V[\beta ,0]+cos(\mu )z_n^S[\beta ,0]+\mathrm{cos}(\frac{\mu }{2})z_n^F[\beta ,0]`$.
We can now make the approximation <sup>4</sup><sup>4</sup>4This approximation can be thought of as a low temperature approximation. This is because at low temperatures $`z_n^S[\beta ,0]`$ approaches zero as $`e^{\beta n}`$. Hence the higher $`z_n`$ are suppressed relative to $`z_1`$. It is also true that for all temperatures, $`z_n<z_1`$ and for very high temperatures we have $`z_n\frac{z_1}{n}`$. As an example, the total contribution for all other $`z_n`$, even near hagedorn transition in free N=4 SYM theory, is only about $`7\%`$ of $`z_1`$. Unitary matrix models involving $`TrU^n,n>1`$ has been discussed in . that $`|z_n[x,i\mu ]|`$ for $`n>1`$ is negligible in comparison to $`|z_1[x,i\mu ]|`$. Neglecting the contribution from the $`n>1`$ modes we arrive at a model which contains only $`TrUTrU^1`$. Using the specific formula for $`z[x,\mu ]`$, of the bosonic <sup>5</sup><sup>5</sup>5Effect of the fermions is discussed in appendix A. sector ,
$`Z(\beta ,Q_0)`$ $`=`$ $`{\displaystyle ๐\mu \mathrm{exp}(i\mu Q_0)DU\mathrm{exp}((a+c\mathrm{cos}(\mu ))TrUTrU^1)}`$ (26)
$`=`$ $`{\displaystyle DU\mathrm{exp}(aTrUTrU^1)๐\mu \mathrm{exp}(i\mu Q_0)\mathrm{exp}(c\mathrm{cos}(\mu )TrUTrU^1)}`$
$`=`$ $`{\displaystyle DU\mathrm{exp}(aTrUTrU^1)I_{Q_0}(cTrUTrU^1)}`$
Here $`a(\beta )=z_n^V[\beta ,0]`$,$`c(\beta )=z_n^S[\beta ,0]`$ and for convenience we did not show the explicit $`\beta `$ dependence in the equations. $`I_n(x)`$ is the Bessel function.
Hence we end up with a matrix model with an effective potential
$$S_{eff}=a(TrUTrU^1)+\mathrm{log}[I_{Q_0}(cTrUTrU^1)]$$
(27)
where $`a>0,c>0`$. We define $`\rho ^2=TrUTrU^1/N^2`$ to get
$$S_{eff}(\rho )=N^2(a\rho ^2+(1/N^2)\mathrm{log}[I_{Q_0}(cN^2\rho ^2)])$$
(28)
It may seem that the logarithmic term is suppressed by a factor of $`1/N^2`$ and hence negligible for large N. But this need not be the case because, in the semicalssical large N limit, we must deal with a system with a charge of order $`N^2`$. Hence we define $`Q_0=N^2q`$ $`(qo(1))`$. Using the asymptotic expansion of $`I_n(nx)`$ for large n, the effective action <sup>6</sup><sup>6</sup>6It should be noted that when Q=0 we should use the asymptotic expansion of $`I_0(x)`$ for large $`x`$. Then we get the expected answer $`S_{eff}=(a+c)TrUTrU^1`$ which is same as a model without any constraint on charge. becomes
$$S_{eff}(\rho )=N^2(a\rho ^2+q(\sqrt{(1+\frac{c^2}{q^2}\rho ^4)}+\mathrm{log}(\frac{\frac{c}{q}\rho ^2}{1+\sqrt{1+\frac{c^2}{q^2}\rho ^4}})))+O(1)$$
(29)
### 3.3 Phase Structure
To understand the phase diagram of this model at large N, we have to locate the saddle points of (29) after including the relevant contribution from the path integral measure depending on whether $`\rho <\frac{1}{2}`$ or $`\rho >\frac{1}{2}`$.<sup>7</sup><sup>7</sup>7The term in the right hand side of equation (31) originates from the path integral measure over an unitary matrix.( see ,Appendix of ).
Differentiating $`S_{eff}(\rho )`$ we get
$`{\displaystyle \frac{}{\rho ^2}}S_{eff}(\rho )`$ $`=`$ $`a+{\displaystyle \frac{}{\rho ^2}}({\displaystyle \frac{1}{N^2}}\mathrm{log}[I_Q(Q{\displaystyle \frac{cN^2\rho ^2}{Q}})]`$ (30)
$`=`$ $`a+b{\displaystyle \frac{I_Q^{}(Q\frac{c\rho ^2}{q})}{I_Q(Q\frac{c\rho ^2}{q})}}`$
$`=`$ $`a+{\displaystyle \frac{q}{\rho ^2}}(1+{\displaystyle \frac{c^2}{q^2}}\rho ^4)^{\frac{1}{2}}+O(1/Q)`$
Hence the equations to solve are
$`a\rho +{\displaystyle \frac{q}{\rho }}(1+{\displaystyle \frac{c^2}{q^2}}\rho ^4)^{\frac{1}{2}}`$ $`=`$ $`\rho ,\rho <{\displaystyle \frac{1}{2}}`$
$`a\rho +{\displaystyle \frac{q}{\rho }}(1+{\displaystyle \frac{c^2}{q^2}}\rho ^4)^{\frac{1}{2}}`$ $`=`$ $`{\displaystyle \frac{1}{4(1\rho )}},\rho >{\displaystyle \frac{1}{2}}`$ (31)
The left side in (31) can be written as
$`a\rho +{\displaystyle \frac{q}{\rho }}(1+{\displaystyle \frac{c^2}{q^2}}\rho ^4)^{\frac{1}{2}}=a\rho +c\rho +{\displaystyle \frac{q}{\rho }}{\displaystyle \frac{1}{(1+\frac{c^2}{q^2}\rho ^4)^{\frac{1}{2}}+\frac{c}{q}\rho ^2}}`$ (32)
So fixing the charge gives rise to a term of type $`\frac{q}{\rho }\frac{1}{(1+\frac{c^2}{q^2}\rho ^4)^{\frac{1}{2}}+\frac{c}{q}\rho ^2}`$ which has some important properties.
1. For all values of $`q>0,c>0`$, this term is a decreasing positive function of $`\rho `$ and it diverges as $`\rho 0`$.
2. For all values of $`c>0`$ it is a monotonically increasing function of $`q`$.
We can now discuss the solution of this model at $`N=\mathrm{}`$. Let us assume that we are discussing the phase where $`a(T)+c(T)<1`$. This condition is valid for low temperatures since $`a(T),c(T)0`$ as $`T0`$. It should also be recalled that without any charge fixing the hagedorn transition occurs when $`a(T)+c(T)=1`$. Unlike the situation with no charge, here we have a function, on the left hand side of (31), which diverges as $`\rho 0`$. Hence $`\rho =0`$ can not be a solution. We get only one solution at a finite value of $`\rho `$ which we will describe in the next paragraph.
Equation (31) is solved in the region $`\rho <\frac{1}{2}`$ with solution
$`\rho ^4={\displaystyle \frac{q^2}{(1a)^2c^2}}={\displaystyle \frac{q^2}{(1ac)(1a+c)}}`$ (33)
The self consistency condition for a solution in the region $`\rho <\frac{1}{2}`$ is
$`\rho ^4={\displaystyle \frac{q^2}{(1a)^2c^2}}<{\displaystyle \frac{1}{16}}`$ (34)
At low temperatures the condition is naturally satisfied for a small enough value of $`q`$. If we gradually increase the temperature (i.e. the value of $`a(T)`$ and $`c(T)`$) while keeping the value of the $`q`$ fixed, then the value of $`\rho `$ at this saddle point will increase. At some temperature $`T_3(q)`$, $`\rho `$ will become equal to $`\frac{1}{2}`$. Since the measure part (i.e. right hand side of (31) ) has a third order discontinuity at $`\rho =\frac{1}{2}`$, we will get a third order phase transition at the temperature $`T_3`$. From (34) we have the following condition at $`T_3`$
$`{\displaystyle \frac{q^2}{(1a)^2c^2}}={\displaystyle \frac{1}{16}}`$ (35)
If the temperature is increased beyond $`T_3`$ then we have to solve (31) in the region $`\rho >\frac{1}{2}`$.
If we increase $`q`$, then the value of $`\rho `$ at the saddle point for a fixed temperature will increase. At some $`q_3`$ we get a third order phase transition satisfying
$`16q_3(T)^2=(a1)^2c^2`$ (36)
Since the minimum value of $`a(T)`$ and $`c(T)`$ is zero, the maximum possible value of $`q_3^2(T)`$ is $`q_{crit}^2=\frac{1}{16}`$. If we increase the $`q`$ beyond $`q_{crit}`$, the saddle point will always be confined in the parameter range $`\rho >\frac{1}{2}`$. Consequently as we increase the temperature from zero we will not get a third order phase transition for $`q>q_{crit}=\frac{1}{4}`$.
This free model, at zero gauge coupling ($`l_s>>R`$ in bulk), has some similarities with $`AdS_5`$ black holes in a fixed charge ensemble. However unlike the three black hole branches in $`AdS_5`$, we get only one branch in the free theory. But most importantly the solution always has a nonzero value of $`\rho `$.
It should be recalled that before the Hagedorn transition, a free gauge theory with zero charge has the solution $`\rho =0`$ .
Some properties of the free theory will be important when we analyze the situation for the weakly coupled gauge theory. Just above the temperature $`T_3(q)`$, the difference of the two sides of (31) can be expanded in the region $`\rho >\frac{1}{2}`$. Defining $`\rho =\frac{1}{2}+x`$,($`x>0`$) the difference is
$`ฯต(q)xC_1x^2`$ (37)
Here $`ฯต(q)>0`$ and $`ฯต(q)0`$ as $`q0`$. It is important to note that $`C_1>2`$ because the measure function (i.e. right hand side of (31) has a third order discontinuity at the point $`\rho =\frac{1}{2}`$. We will also discuss the significance of this in what follows.
## 4 Small coupling model
We will now discuss the problem with a small non-zero gauge coupling $`\lambda =g_{YM}^2N`$. By AdS/CFT correspondence it corresponds to a finite string length in $`AdS`$. It has been shown in that by considering a phenomenological model of type
$`S[TrUTrU^1]=a(\lambda ,T)(TrUTrU^1)+{\displaystyle \frac{b(\lambda ,T)}{N^2}}(TrUTrU^1)^2,b>0`$ (38)
we can map out the possible phase diagram of type $`IIB`$ string theory in $`AdS_5`$.<sup>8</sup><sup>8</sup>8In fact in an arbitrary convex function is considered, and shown to map out the phase diagram of $`IIB`$ theory in $`AdS_5`$. The simplified model (38) leads to similar qualitative result. Even though the model in (38) can be derived from a weak coupling analysis of the gauge theory, it can be thought of as a phenomenological model describing supergravity in $`AdS_5`$.
We are motivated to discuss the fixed charge ensemble in the same spirit. Let us add a small interaction term to (28). The effective action is then given by
$`S_{eff}(\rho )=N^2(a\rho ^2+b\rho ^4+q(\sqrt{(1+{\displaystyle \frac{c^2}{q^2}}\rho ^4)}+\mathrm{log}({\displaystyle \frac{\frac{c}{q}\rho ^2}{1+\sqrt{1+\frac{c^2}{q^2}\rho ^4}}})))+O(1)`$ (39)
Here $`b`$ is proportional to $`\lambda `$ and is also a function of charge. Depending on the theory considered, the sign of $`b`$ can be either positive or negative. It has been shown in that $`b`$ is positive for a pure YM theory. In the following discussion we will assume that this is the case in order to motivate a similarity with the supergravity picture. The equations determining the saddle points of (39) , including the contribution from the path integral measure, are
$`(a+c)\rho ^2+2b\rho ^4+{\displaystyle \frac{q}{(1+\frac{c^2}{q^2}\rho ^4)^{\frac{1}{2}}+\frac{c}{q}\rho ^2}}`$ $`=`$ $`\rho ^2,\rho <{\displaystyle \frac{1}{2}}`$
$`(a+c)\rho ^2+2b\rho ^4+{\displaystyle \frac{q}{(1+\frac{c^2}{q^2}\rho ^4)^{\frac{1}{2}}+\frac{c}{q}\rho ^2}}`$ $`=`$ $`{\displaystyle \frac{\rho }{4(1\rho )}},\rho >{\displaystyle \frac{1}{2}}`$ (40)
In what follows it will be useful to introduce a function
$`M(\rho )`$ $`=`$ $`\rho ^2,\rho <{\displaystyle \frac{1}{2}}`$ (41)
$`=`$ $`{\displaystyle \frac{\rho }{4(1\rho )}},\rho >{\displaystyle \frac{1}{2}}`$
$`M(\rho )`$ is an increasing convex function and is the right hand side of (40). It is also useful to introduce $`D(x)=S_{eff}^{}(x)M(x)`$. Eqn (40) is then equivalent to $`D(\rho )=0`$.
It has been shown in that for an interacting model with zero charge, we get nucleation of black holes along the curve given by $`a=\frac{12w}{(1w)^2(1+w)}`$ and $`b=\frac{2w}{(1w)^2(1+w)^3}`$, $`1>w>0`$. Here we want to analyze a similar type of phenomenon. Let us consider the different cases.
### $`T`$ is varying and $`q`$ is small
Let us start with a value of charge which is small (i.e. $`q1`$) and let us increase the temperature from zero. At low temperatures all the parameters $`a,c`$ will be small (for small $`T`$ these parameters have a dependence like $`e^{c\beta }`$, $`c`$ is a constant ). Hence we get just one solution for $`\rho <\frac{1}{2}`$ which we call $`I`$. There is no solution for $`\rho >\frac{1}{2}`$ (see the topmost curve of Fig 2) because the left hand side of (40) is less than the right hand side.
The situation is quite similar to supergravity where for small charge and low enough temperatures we get a stable small black hole solution. <sup>9</sup><sup>9</sup>9We should keep in mind that $`\rho <\frac{1}{2}`$ is not the supergravity regime. The solutions of gauge theory effective action there should be represented as excited string states.
The function $`M(\rho `$) (i.e. right hand side of (40)) is a convex increasing function. Hence we will generate new solutions of (40) in the region $`\rho >\frac{1}{2}`$ as we increase temperature (i.e. $`a(T),c(T)`$ as discussed in appendix B) keeping $`q`$ fixed (Fig 2, Fig 3). The new solutions will always come in pairs (Fig 2). Let us call the solution nucleation temperature as $`T_{01}(q)`$. At $`T=T_{01}`$ we have
$`D(\rho )=S_{eff}^{}(\rho )M(\rho )`$ $`=`$ $`0`$
$`D^{}(\rho )=0`$ (42)
As we increase the temperature further (i.e. $`a(T),c(T)`$) the function on the left hand side of (40) will also increase. Hence the solutions will start to separate. Let us call them $`II`$ and $`III`$. Here $`\rho _{III}>\rho _{II}`$. Also, $`III`$ is a stable saddle point whereas $`II`$ is an unstable one. They are similar to stable big and unstable small black holes in supergravity. As the temperature is increased beyond $`T_{01}`$ the value $`\rho _{II}`$ decreases whereas $`\rho _{III}`$ increases. At some temperature $`T_H(q)`$, we will have $`S(III)<S(I)`$ and consequently we expect a first order phase transition. At a temperature $`T_{HP}`$, the dominant saddle point of the system changes from $`I`$ to $`III`$. As the temperature increases the saddle point $`II`$ goes through a third order transition when $`\rho (II)`$ crosses the $`\rho =\frac{1}{2}`$ point. Call this temperature $`T_3`$ which is determined by the following relation between the parameters
$`(a+{\displaystyle \frac{b}{2}})+4q(1+{\displaystyle \frac{c^2}{16q^2}})^{\frac{1}{2}}=1`$ (43)
Increasing the temperature further, the saddle point $`II`$ approaches the saddle point $`I`$ and they merge at a temperature $`T_{02}`$ ($`4th`$ graph from above in Fig 2 ). In the language of catastrophe theory this is a fold catastrophe. For $`T>T_{02}`$, the only saddle point is $`III`$. This then is the thermal history as we increase the temperature for a small $`q`$.
In summary at low temperatures, we have only one saddle point $`I`$ and then two new saddle point $`II,III`$ are created at $`T_{01}`$. As the temperature increases the saddle points $`I,II`$ merge at a temperature $`T_{02}`$. Beyond that we have only one saddle point $`III`$. In the next paragraph we will discuss what happens when we increase the value of the $`q`$.
### Varying q$`:`$
Let us discuss how the various temperatures discussed above change as we increase the value of $`q`$.
1. The first is $`T_{01}(q)`$, the nucleation temperature for saddle points $`II`$ and $`III`$. $`T_{01}`$ will decrease as we increase the value of $`q`$. This is so because all three terms in the right hand side of (40) are positive and increasing functions of $`\rho `$ and the left hand side is a positive convex function. Consequently the saddle point value of $`\rho `$ at $`T_{01}`$ will also decrease.
2. The temperature $`T_{02}`$ at which the stable saddle point $`I`$ and unstable saddle point $`II`$ merge and the value of $`\rho `$ at $`T_{02}`$, will increase as we increase $`q`$. The reason is that all the coefficients in (40) are positive and hence increasing the value of $`q`$ increases the function on the left hand side.
As we increase the value of $`q`$ further, $`T_{01}`$ and $`T_{02}`$ will become equal for some value of $`q=q_{crit}`$. In the language of catastrophe theory this is a cusp catastrophe. Corresponding to $`q_{crit}`$ there will be a $`T_{crit}=T_{01}=T_{02}`$ and a value of $`\rho _{crit}`$ (saddle point $`\rho `$ at $`q_{crit}`$ , $`T_{crit}`$). As we increase the $`q`$ beyond $`q_{crit}`$, we do not get any new saddle point and consequently there is only one saddle point for all temperatures. We will discuss the physics near this phase transition in detail in what follows.
### Value of $`\rho _{crit}`$ :
Let us determine the value of $`\rho _{crit}`$. As already discussed at the end of the previous chapter, in (37) $`C`$ is always a finite quantity. Hence, (40) cannot have three solutions in the region $`\rho <\frac{1}{2}`$ for small $`b`$. Therefore the saddle point $`I`$ is always in the region $`\rho <\frac{1}{2}`$. Whereas the saddle point $`III`$ will be in the region $`\rho >\frac{1}{2}`$. Hence the only place where these three saddle points can meet is $`\rho _{crit}=\frac{1}{2}`$ which is also the point of the third order phase transition.
### Physical Interpretation
Before proceeding further we will briefly discuss the bulk interpretation of saddle points of weakly coupled gauge theory. Weak coupling in gauge theory means $`l_sR_{AdS_5}`$. Hence the supergravity picture is not valid in the bulk. However at large $`N`$, the string coupling (i.e. $`\frac{1}{N}`$) will be small and we may conclude that the saddle points discussed above can be described by exact (in all orders in $`l_s`$) conformal field theories. These CFTs are characterized by the values of $`q`$ and $`\rho `$ at the saddle points.
We would like to end this section by emphasizing that the coincidence of the three saddle points at $`\rho =\frac{1}{2}`$, is a property of the weak coupling ($`\lambda <<1`$) limit of the gauge theory. In what follows we will see that this fact is not necessarily true at strong coupling. We will show that there the coincidence happens in the gaped phase where $`\rho >\frac{1}{2}`$.
## 5 Effective action and phase diagram at strong coupling
In this section we will discuss the effective action and the phase diagram in the strongly coupled gauge theory which is dual to the supergravity (discussed in section 2) regime of $`IIB`$ string theory.
### 5.1 Finite temperature effective actions in the gauge theory
Let us first summarize the situation in the zero charge sector. The propagator of adjoint fields in the free gauge theory, on a compact manifold $`S^3`$, coupled to a space independent $`A_0`$ is given by (see )
$`G_{Ukl}^{ij}(x,t,y,0)={\displaystyle \underset{n=\mathrm{}}{\overset{n=\mathrm{}}{}}}(U^n)_j^iG_0(x,t+n\beta ,y,0)(U^n)_j^k`$ (44)
where $`G_0(x,t,y,0)`$ is the zero temperature Greenโs function and $`U`$ is the constant Polykov line. We know that at any temperature <sup>10</sup><sup>10</sup>10Temperature here is measured in units of $`\frac{1}{R_{S^3}}`$. $`G_0(x,t+n\beta ,y,0)>G_0(x,t,y,0)`$ and also at low temperatures
$`G_0(y,t+n\beta ,x,0)e^{n\beta }`$ (45)
Using the above Greenโs function one can develop the large N diagrammatic to arrive at an effective action involving $`Z_N`$ invariant terms built out of products of $`trU^n`$. In fact one can imagine integrating out all the modes $`trU^n`$ for $`n>1`$ in favor of $`trUtrU^1`$. In this way one gets an effective action of the form
$`S_{eff}={\displaystyle \underset{n=1}{\overset{n=\mathrm{}}{}}}a_n(\beta ,\lambda )({\displaystyle \frac{TrUTrU}{N^2}})^n`$ (46)
As one increases the coupling constant we would expect that the form of the effective action would remain the same except that the dependence of the parameters on temperature and the โthooft coupling would change.
### 5.2 Non-zero charge sector
If we include the fixed charge constraint, as in (25), then we get the following expression for the fixed charge path integral
$`Z_q={\displaystyle ๐\mu e^{i\mu Q}DUe^{N^2(_{n=0}^{n=\mathrm{}}S_n(\rho ,\lambda ,\beta )\mathrm{cos}(n\mu ))}}`$ (47)
For large $`N`$ we can do the $`\mu `$ integral by the saddle point method. The saddle point of $`\mu `$ is on the imaginary axis. Hence we set $`\mu =im`$, to get the saddle point equation
$`q={\displaystyle \underset{n=1}{\overset{n=\mathrm{}}{}}}nS_n(\rho ,\lambda ,\beta )\mathrm{sinh}(nm)),q=Q/N^2`$ (48)
At small values of $`\rho `$, $`S_n(\rho )`$ goes as $`\rho ^n`$ <sup>11</sup><sup>11</sup>11Introduction of a chemical potential changes the formula (44) as
$$G_{Ukl}^{\mu ij}(x,t,y,0)=e^{\frac{i\mu t}{\beta }}\underset{n=\mathrm{}}{\overset{n=\mathrm{}}{}}(U^n)_j^iG_0^\mu (x,t+n\beta ,y,0)(U^n)_j^k$$
where
$$G_0^\mu (x,t+n\beta ,y,0)=e^{in\mu }G_0(x,t+n\beta ,y,0)$$
(49) Hence in each order in perturbation theory the terms containing $`\mathrm{cos}(n\mu )`$ also get multiplied by $`(TrUTrU^1)^m,m>n`$ or the higher operators like $`TrU^n`$ which can be integrated out to give again a term like $`(TrUTrU^1)^m`$. and hence in the $`\rho 0`$ limit we can approximate the equation (48) as
$`q๐\rho \mathrm{sinh}(m)`$ (50)
where $`๐`$ is a constant independent of $`\rho `$. The solution is
$`mq\mathrm{log}(\rho )+๐,\rho 0`$ (51)
Substituting $`m`$ in (47) we get a logarithmic term for $`\rho `$. We conclude that the logarithmic term is a general feature and it implies among other things that $`TrU=0`$ is never a solution in the non-zero charge sector.
### 5.3 Model effective action at strong coupling
Following our previous discussion, we will include the generic logarithmic term in the effective potential for a fixed non-zero charge, and proceed to analyze the saddle point structure following .
Our proposal for the gauge theory effective action is
$`S_q=S(a(T),b(T),\mathrm{},\rho )+q\mathrm{log}(\rho )`$ (52)
and the saddle point equations are
$`\rho F(a(T),b(T),c(T),\mathrm{}\rho ^2)+q=\rho ^2,\rho <{\displaystyle \frac{1}{2}}`$
$`\rho F(a(T),b(T),c(T),\mathrm{}\rho ^2)+q={\displaystyle \frac{\rho }{4(1\rho )}},\rho >{\displaystyle \frac{1}{2}}`$ (53)
Where $`F(\rho )=S_{eff}^{}(\rho )`$. We assume that
1. $`F(x,T)`$ is a monotonically increasing function of $`x`$ and F(0,T)=0
2. Value of $`F(x,T)`$ increases for fixed $`x`$ as we increase the temperature and F(x,0)=0.
These global properties of $`F(x)`$ reproduce the phase diagram of supergravity.
### Analysis of solution structure
Let us consider the function $`D(T,\rho )=\rho F(\rho ,T)M(\rho )`$ (Where $`M(\rho )`$ is the contribution from measure appearing at the right hand side of (53)). At $`T=0`$, $`F(\rho ,T)`$ is zero. Hence $`D(T,\rho )`$ is a monotonically decreasing function of $`\rho `$ at $`T=0`$.
We know that at $`T=T_{01}`$ a pair of two new saddle points appear at $`\rho _{01}>\frac{1}{2}`$. Hence at $`T=T_{01}`$ we have $`D(T_{01},\rho _{01})=0`$ and $`D^{}(T_{01},\rho _{01})=0`$. $`D(T_{01},\rho )`$ has a zero for $`\rho =0`$ and it is a decreasing function in the neighborhood of $`\rho =0`$. It again increases and become zero at $`\rho _{01}`$ and then the function again decreases as $`\rho 1`$. This implies that the function has a local maximum and local minimum.
In summary:
1. $`D(0,\rho )`$ is a monotonically decreasing function of $`\rho `$.
2. $`D(T_{01},\rho )`$ has a maximum and minimum.
There is a temperature $`T_{crit}`$ at which the local maximum and local minimum appear (Fig 4). Let us call this temperature $`T_{crit}`$. At $`T_{crit}`$ the curve $`D(T_{crit},\rho )`$ will have a point of inflection at $`\rho =\rho _{crit}`$, say. Let the value of $`D(T_{crit},\rho _{crit})=q_{crit}`$.
Increasing the value of $`q`$ from zero we need to solve the equation $`D(T,\rho )=q`$. We will get a solution for a non-zero value of $`\rho `$. Denote this solution as $`I`$. As the temperature increases, two new solutions appear at $`T=T_{01}`$. Call the stable solution as $`III`$, and the unstable solution $`II`$. As the temperature is further increased to $`T_{02}`$, the unstable solution $`II`$ and the stable solution $`I`$ merge. For $`T>T_{02}`$ , the only solution is $`III`$.
As $`q`$ approaches $`q_{crit}`$ from below the two temperatures $`T_{01}`$ and $`T_{02}`$ approach each other. At $`q=q_{crit}`$, we have $`T_{01}=T_{02}=T_{crit}`$. If we increase $`q`$ beyond $`q_{crit}`$ only one solution appears for all temperature. These facts are consistent with supergravity solutions(section 2).
With a sufficiently sharp rising function $`F(T,\rho )`$ in (53) we can obtain this critical point in the region $`\rho >\frac{1}{2}`$.<sup>12</sup><sup>12</sup>12Which describes supergravity in the bulk.. As the function $`D(\rho ,T)`$ is smooth in the region $`\rho >\frac{1}{2}`$ the second derivative of the function $`D(T,\rho )`$ will vanish at the inflection point, and we will get a third order phase transition. We can calculate the partition function in suitable double scaling limit near the critical point. This is discussed in section 6.
### A specific example
We will now illustrate the above phenomenon in a simple model defined by
$`F(\rho )=a\rho +b\rho ^3`$ (54)
where $`a,b>0`$.
We will determine the parameter ranges of $`a,b`$ for which all the three saddle points of (54) are in the range $`\rho >\frac{1}{2}`$.
At $`\rho =\frac{1}{2}`$ we have the constraints
$`_\rho (\rho F(\rho ){\displaystyle \frac{\rho }{4(1\rho )}})<0`$
$`_\rho ^2(\rho F(\rho ){\displaystyle \frac{\rho }{4(1\rho )}}>0`$ (55)
Putting the value of $`\rho =\frac{1}{2}`$ in the above inequality we get the following constraints on the parameters $`a+b<1`$ and $`a+3b>2`$. Simplifying we have $`b>\frac{1}{2}`$ and $`a<\frac{1}{2}`$.
As the coupling become stronger, we expect that $`b`$ is not necessarily small and will be of $`o(1)`$ or greater. All the saddle points of (53) are then naturally shifted to the region $`\rho >\frac{1}{2}`$. Here, as was discussed in , we can expect to match the solutions of the gauge theory with those of supergravity. The stable saddle point $`I`$ corresponds to the stable black hole branch $`I`$ of supergravity. And unstable saddle point $`II`$ is matched with the unstable black hole branch $`II`$. The stable saddle point $`III`$ is matched with the big stable black hole in supergravity. With this identification the thermal history and critical behavior of the gauge theory, discussed earlier in this chapter, match with the thermal history and critical behavior of supergravity (discussed in section 2 and ).
## 6 Universal neighborhood of critical point and the critical exponents
Let us consider the effective action $`S_{tot}(\rho ,T,q)`$ which includes the contribution from the path integral measure over an unitary matrix. The derivative of $`S_{tot}`$ with respect to $`\rho `$, say $`G(\rho ,T,q)`$, gives the saddle point equations (53). We have already discussed that by a suitable choice of parameters the critical point appears in the region $`\rho >\frac{1}{2}`$. This critical point is a third order critical point because three saddle points of the system merges here. Hence the first and second derivatives of $`G_{tot}(\rho ,T,q)`$ with respect to $`\rho `$ vanish at $`\rho =\rho _{crit},q=q_{crit},T=T_{crit}`$. Expanding $`G(\rho ,T,q)`$ around the critical point , we get
$$G(\rho ,T,q)=(\delta \rho )^3\frac{_\rho ^3G}{3!}+(\delta T)_TG+(\delta q)_qG+(\delta \rho )(\delta T)_\rho _TG+(\delta q)(\delta T)_\rho _qG$$
(56)
Let us fix $`T=T_{crit}`$ or $`\delta T=0`$. Then the equation (56) has one solution. In order to know how the saddle point value of $`\rho `$ approaches $`\rho _{crit}`$ ($`\delta \rho 0`$) as $`\delta q0`$, we equate the leading part of (56) to zero.
$`(\delta \rho )^3{\displaystyle \frac{_\rho ^3G}{3!}}+(\delta q)_qG=0`$ (57)
Hence $`\delta \rho \delta q^{\frac{1}{3}}`$ and we get the same universal exponent $`\frac{1}{3}`$ as in supergravity .
### 6.1 Partition function near the critical point
Near the critical point we can write the $`S_{tot}`$ as
$`S_{tot}=S_{tot}(\rho _{crit},T_{crit},q_{crit})+(\delta \rho )^4{\displaystyle \frac{_\rho ^4S}{4!}}+(\delta q)_qS+(\delta q)(\delta \rho )_\rho _qS+O(\delta \rho ^5)`$ (58)
If we define a double scaling limit $`N^{\frac{1}{2}}\rho =x,N^{\frac{3}{2}}q=z`$ , we can write the $`o(1)`$ part of the partition function, after suitable rescaling of the variables, as
$`Z_2`$ $`{\displaystyle ๐xe^{(x^4zx)}}`$ (59)
This can be calculated in a power series
$`Z_2{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{z^{2n}}{(2n)!}}\mathrm{\Gamma }({\displaystyle \frac{n}{2}}+{\displaystyle \frac{1}{4}})`$ (60)
### 6.2 Approaching the critical point through a line of first order transitions
Another type of double scaling limit is possible in this problem. We can set<sup>13</sup><sup>13</sup>13It is same as following the HP(first order) transition line in parameter space.
$`(\delta T)_TG+(\delta q)_qG=0`$ (61)
by choosing a suitable relation between $`\delta T`$ and $`\delta q`$. Using (61) in (56) we get
$`(\delta \rho )^3{\displaystyle \frac{_\rho ^3G}{3!}}+(\delta \rho )((\delta T)_\rho _TG+(\delta q)_\rho _qG)=0`$ (62)
with the solutions
$`\delta \rho =0`$
$`\delta \rho \pm (\delta T)^{\frac{1}{2}}`$ (63)
We can expand $`S_{eff}`$ as
$`S_{eff}S_{crit}+(\delta \rho )^4{\displaystyle \frac{_\rho ^4S}{4!}}+C_1(\delta T)(\delta \rho )^2+OT`$ (64)
where $`OT`$ are terms independent of $`\delta \rho `$. Defining a suitable double scaling limit. $`N^{\frac{3}{2}}\delta T=z,N^{\frac{1}{2}}\delta \rho =x`$ and a suitable rescaling of the parameters we can evaluate the $`o(1)`$ factors in the partition function as
$`Z_2{\displaystyle ๐xe^{(x^4+2zx^2)}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(2z)^n}{n!}}\mathrm{\Gamma }({\displaystyle \frac{n}{2}}+{\displaystyle \frac{1}{4}})\sqrt{z}e^{\frac{z^2}{2}}K_{\frac{1}{4}}({\displaystyle \frac{z^2}{2}})`$ (65)
where $`z<0`$.
## 7 Conclusion
In this paper we have studied the logarithmic matrix model generated by fixing the R-charge in the gauge theory partition function. In the free gauge theory it has been shown that there is no solution with $`\rho =0`$ ($`AdS`$ type solution). We then studied the effect of adding an interaction term in our model and discussed the generic nature of the logarithmic term even at arbitrary value of the coupling. We identified the supergravity saddle points and their critical behavior which was discussed in ().
Our main aim was to give another example of the utility of unitary matrix methods in providing a non-perturbative dual description of blakholes in $`AdS`$ and to understand the relation between matrix models and string theory in general. It would be interesting to consider an effective unitary matrix model to describe phases of Kerr-Ads black holes.
## 8 Acknowledgment
We would like to thank the theory division of CERN for hospitality where part of this work was done. We acknowledge useful discussions with Luis Alvarez-Gaume and Marcos Marino on many aspects of the matrix model/string theory correspondence. We also acknowledge a correspondence with Hong Liu. PB likes to acknowledge CSIR for SPM fellowship and TIFR alumni association for partial financial support. PB also likes to thank Swagato Mukherjee for technical help in preparation of this paper.
## Appendix A Appendix: Inclusion of Fermions
Including the contributions from the fermions of $`N=4`$ $`SYM`$ theory change (26) to
$`Z(\beta ,Q_0)={\displaystyle DU๐\mu \mathrm{exp}(N^2(a+c\mathrm{cos}(\mu )+d\mathrm{cos}(\frac{\mu }{2}))\rho ^2i\mu Q_0)}`$ (66)
Where $`d(\beta )`$ is the single particle partition function for the fermions.
At large $`N`$ , the integral in (66) could be evaluated by the saddle point method. The equations determining the saddle points of $`\mu =im`$ and $`\rho `$ are
$$(c\mathrm{sinh}(m)+\frac{d}{2}\mathrm{sinh}(\frac{m}{2}))=\frac{q}{\rho ^2}$$
(67)
and
$$\rho (a+\mathrm{cosh}(m)+d\mathrm{cosh}(\frac{m}{2}))=\rho $$
(68)
We would like to see weather there is a solution with $`\rho =0`$. As the right hand side of (67) becomes large in the limit $`\rho 0`$, we can self consistently approximate $`\mathrm{cosh}(m)`$ and $`\mathrm{sinh}(m)`$ as $`e^m`$ and we get
$$m\mathrm{log}\frac{q}{c\rho ^2}$$
(69)
Hence a logarithmic potential for $`\rho `$ is once again generated. One can also confirm this by putting (69) in (68).
## Appendix B Appendix: Positivity of the coefficient of the quadratic term in the effective action
Let us consider the partition function of YM theory on a compact manifold written as an integral over the effective action of $`\rho =TrUTrU^1`$.
$`Z(\beta )={\displaystyle DUe^{N^2(S_{eff}(\rho ))}}`$ (70)
$`={\displaystyle ๐\rho e^{N^2(S_{eff}(\rho )S_M(\rho ))}}`$ (71)
Where $`S_M(\rho )`$ is the contribution from the measure part<sup>14</sup><sup>14</sup>14See discussions before (30). of path integral and
$`S_{eff}(\rho )=a(\beta )\rho ^2+{\displaystyle \underset{n=4}{}}a_n(\beta )\rho ^n`$ (72)
i.e. a polynomial in $`\rho `$. As $`\beta \mathrm{}`$ we have $`S_{eff}(\rho )0`$. Contribution from the measure part $`S_M(0)`$ has only one minimum at $`\rho =0`$. Hence at low temperature the system will have a saddle point at $`\rho =0`$. Expanding $`\rho `$ around this saddle point as $`\rho =0+\frac{\delta \rho }{N}`$ we get
$`Z(\beta )={\displaystyle _{\mathrm{}}^{\mathrm{}}}d(\delta \rho )e^{(1a(\beta ))(\delta \rho )^2+o(\frac{1}{N^2})}`$ (73)
Like any thermal partition function, (71) or (73) is a decreasing function of the $`\beta `$. Hence $`a(\beta )`$ should also be a decreasing function of $`\beta `$. Since $`a(\mathrm{})=0`$, for any finite $`\beta `$, $`a(\beta )`$ is a positive decreasing function. Hagedorn transition happens when $`a(\beta _H)=1`$, but whether $`a(\beta )`$ will reach $`1`$ or not depends on the model.
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# Simple decompositions of simple special Jordan superalgebras
## 1 Introduction
In our previous work with T.Tvalavadze we considered special simple finite-dimensional Jordan algebras decomposable as the sum of two proper simple subalgebras. The main result in is the following.
Theorem. Let $`๐ฅ`$ be a finite-dimensional special simple Jordan algebra over an algebraically closed field $`F`$ of characteristic not two. The only possible decompositions of $`๐ฅ`$ as the sum of two simple subalgebras $`๐ฅ_1`$ and $`๐ฅ_2`$ are the following:
1. $`๐ฅFV`$ and $`๐ฅ_1FV_1`$, $`๐ฅ_2FV_2`$, where $`V`$, $`V_1`$, $`V_2`$ are vector
spaces.
2. Either $`๐ฅH(_3)`$ and $`๐ฅ_1H(F_3)`$, $`๐ฅ_2FV`$, or $`๐ฅH(_n)`$,
$`n3`$, $`๐ฅ_1H(F_n)`$ and $`๐ฅ_2`$ is isomorphic to one of the following
algebras: $`H(F_{n1})`$, $`H(F_n)`$ or $`H(_{n1})`$.
3. $`๐ฅH(๐ฌ_n)`$ and $`๐ฅ_1`$, $`๐ฅ_2H(_n)`$.
Actually the problem of simple decompositions of simple algebras first arises in the paper of Onishchik (see ) in which he classified all possible types of simple decompositions of simple complex and real Lie algebras. Later for associative algebras over an arbitrary field $`F`$ the same problem was formulated and then solved by Bahturin and Kegel in . According to , no full-matrix algebra can be written as the sum of two full-matrix subalgebras. Note that if $`F`$ is algebraically closed with zero characteristic, then this follows from .
To begin with we briefly remind the classification of simple Jordan superalgebras obtained by Kac (see ) over an algebraically closed field $`F`$ with zero characteristic. If $`๐ฅ`$ is a simple special finite-dimensional Jordan superalgebra over algebraically closed field $`F`$ with zero characteristic, then $`๐ฅ`$ is isomorphic to one of the following superalgebras:
(1) $`M_{n,m}(F)^{(+)}`$, the set of all matrices of order $`n+m`$ with respect to the natural $`Z_2`$-gradation under the Jordan supermultiplication;
(2) $`osp(n,m)`$, the set of all matrices of order $`n+2m`$ symmetric with respect to the orthosymplectic superinvolution. The superalgebra consists of matrices $`\left\{\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)\right\}`$ where $`A^t=A`$, $`D`$ is symplectic, $`B`$, $`C`$ are skew-symmetric;
(3) $`P(n)=\{\left(\begin{array}{cc}A& B\\ C& A^t\end{array}\right),B^t=B,C^t=CM_n(F)\}`$;
(4) $`Q(n)=\left\{\left(\begin{array}{cc}A& B\\ B& A\end{array}\right)\right\}`$ where $`A`$ and $`B`$ are any square matrices of order $`n`$.
(5) Let $`V=V_0+V_1`$ be a $`Z_2`$-graded vector space with a non-singular symmetric bilinear superform $`f(,):V\times VF`$. Consider the direct sum of $`F`$ and $`V`$, $`๐ฅ=FV`$, and determine multiplication according to
$$(\alpha +v)(\beta +w)=(\alpha \beta +f(v,w))+(\alpha w+\beta v).$$
Then $`๐ฅ`$ becomes a Jordan superalgebra of the type $`J(V,f)`$ with respect to the following $`Z_2`$-gradation: $`๐ฅ_0=F+V_0`$, $`๐ฅ_1=V_1`$.
(6) The 3-dimensional Kaplansky superalgebra $`K_3`$, $`(K_3)_0=Fe`$, $`(K_3)_1=Fx+Fy`$, with the multiplication $`e^2=e`$, $`ex=\frac{1}{2}x`$, $`ey=\frac{1}{2}y`$, $`[x,y]=e`$.
(7) The 1-parametric family of 4-dimensional superalgebras $`D_t`$, $`D_t=(D_t)_0+(D_t)_1`$, where $`(D_t)_0=Fe_1+Fe_2`$, $`(D_t)_1=Fx+Fy`$, where $`e_i^2=e_i`$, $`e_1e_2=0`$, $`e_ix=\frac{1}{2}x`$, $`e_iy=\frac{1}{2}y`$, $`[x,y]=e_1+te_2`$, $`i=1,2`$. A superalgebra $`D_t`$ is simple only if $`t0`$. If $`t=1`$, then $`D_1`$ isomorphic $`M_{1,1}(F)`$.
Next we cite some important Lemmas and Theorems from which will be repeatedly used later.
###### Lemma 1.1
Let a Jordan algebra $`๐ฅ`$ of the type $`H(๐_m^{})`$ be a proper subalgebra of $`H(๐_n)`$ such that the identity of $`H(๐_n)`$ is an element of this subalgebra. If either
1. $`๐^{}=F`$ and $`๐=F`$, or
2. $`๐^{}=,๐ฌ`$, and $`๐=`$,
then $`m\frac{n}{2}`$
###### Lemma 1.2
Let $`V`$ be a vector space with a non-singular symmetric bilinear form $`f`$, and $`v_0`$ a fixed non-trivial vector in $`V`$. Let $`๐ฎ`$ be the set of all linear operators which are symmetric with respect to $`f`$. Then, $`๐ฎv_0=V`$.
###### Theorem 1.3
Let $`๐ฅ`$ be a simple Jordan algebra of the type $`H(F_n)`$, and $`๐`$, $``$ proper simple Jordan subalgebras of $`๐ฅ`$. Then $`๐ฅ๐+`$.
###### Lemma 1.4
Let $`W`$ be the natural $`2m`$-dimensional module for $`H(๐ฌ_m)`$, and $`v`$ be an arbitrary non-zero vector in $`W`$. Then, $`dimH(๐ฌ_m)v=2m1`$.
Throughout the paper the basic field $`F`$ is algebraically closed with characteristic zero.
## 2 Decompositions of superalgebras of the type $`M_{n,m}(F)^{(+)}`$
Our main goal is to prove the following.
###### Theorem 2.1
Let $`๐`$ be a superalgebra of the type $`M_{n,m}(F)^{(+)}`$ where $`n,m>0`$. If both $`n,m`$ are odd, then $`๐`$ has no decompositions into the sum of two proper nontrivial simple subsuperalgebras. If one of the indices, for example $`m`$, is an even number and the other is odd, then the only possible simple decomposition is the following: $`๐=+๐`$ where $``$ and $`๐`$ have types $`osp(n,\frac{m}{2})`$ and $`M_{n1,m}(F)^{(+)}`$, respectively. If both indices are even, then $`๐`$ admits two types of decompositions of the following forms:
1. $`๐=_1+๐_1`$ where $`_1`$ and $`๐_1`$ have types $`osp(n,\frac{m}{2})`$ and $`M_{n1,m}(F)^{(+)}`$,
2. $`๐=_2+๐_2`$ where $`_2`$ and $`๐_2`$ have types $`osp(m,\frac{n}{2})`$ and $`M_{m1,n}(F)^{(+)}`$.
Before the discussion of various properties of $`M_{n,m}(F)^{(+)}`$ we recall a definition of the universal associative enveloping superalgebra of a Jordan superalgebra which will be frequently used later.
An associative specialization $`u:๐ฅU(๐ฅ)`$ where $`U(๐ฅ)`$ is an associative superalgebra is said to be universal if $`U(๐ฅ)`$ is generated by $`u(๐ฅ)`$, and for any other specialization $`\phi :๐ฅ๐`$ where $`๐`$ is an associative superalgebra there exists a homomorphism $`\psi :U(๐ฅ)๐`$ such that $`\phi =\psi u`$. Then $`U(๐ฅ)`$ is called a universal associative enveloping superalgebra of $`๐ฅ`$. It is worth noting that an associative superalgebra can be considered as an associative algebra. The following Theorem by C.Martinez and E.Zelmanov plays a key role in the later discussion.
###### Theorem 2.2
Let $`U(๐ฅ)`$ denote a universal associative enveloping superalgebra for a Jordan superalgebra $`๐ฅ`$. Then $`U(M_{k,l}^{(+)})M_{k,l}(F)M_{k,l}(F)`$ where $`(k,l)(1,1)`$; $`U(Q(k))Q(k)Q(k)`$, $`k2`$; $`U(osp(m,n))M_{m,2n}(F)`$, $`(m,n)(1,1)`$; $`U(P(n))M_{n,n}(F)`$, $`n3`$.
Remark 1(see ) In the case where $`๐ฅM_{1,1}(F)^{(+)},P(2),osp(1,1),K_3`$ or $`D_t`$ the universal enveloping superalgebras have more complicated structure. Indeed, the universal associative enveloping superalgebras of the above Jordan superalgebras are no more finite-dimensional. Also we note that if the characteristic of the basic field $`F`$ equals zero, then $`K_3`$ has no non-zero finite-dimensional associative specializations.
The following Theorem by Martinez and Zelmanov (see ) describes all irreducible one-sided bimodules of $`D(t)`$ where $`t1,0,1`$.
###### Theorem 2.3
Let F be an algebraically closed field with zero characteristic. If $`t=\frac{m}{m+1}`$, $`m1`$, then $`D(t)`$ has two irreducible finite-dimensional one sided bimodules $`V_1(t)`$ and $`V_1(t)^{op}`$.
If $`t=\frac{m+1}{m}`$, $`m1`$, then $`D(t)`$ has two irreducible finite-dimensional one sided bimodules $`V_2(t)`$ and $`V_2(t)^{op}`$.
If $`t`$ cannot be represented as $`\frac{m}{m+1}`$ and $`\frac{m+1}{m}`$ where $`m`$ is a positive integer, then $`D(t)`$ does not have non-zero finite-dimensional specializations.
Remark 2 If $`\text{char}F=p>2`$, then for an arbitrary $`t`$ the superalgebra $`D(t)`$ can be embedded in a finite-dimensional associative superalgebra.
Remark 3 If $`t=\frac{m}{m+1}`$, then $`dimV_1(t)_0=m`$, $`dimV_1(t)_1=m+1`$. If $`t=\frac{m+1}{m}`$, then $`dimV_2(t)_0=m+1`$, $`dimV_2(t)_1=m`$.
Now we look at the case when $`๐ฅJ(V,f)`$. Let $`V=V_0+V_1`$ be a $`Z_2`$-graded vector space, $`dimV_0=m`$, $`dimV_1=2n`$. Let $`f(,):V\times VF`$ be a supersymmetric bilinear form on $`V`$. The universal associative enveloping algebra of the Jordan algebra $`F+V_0`$ is the Clifford algebra $`C(V_0,f)=1,e_1,\mathrm{},e_m|e_ie_j+e_je_i=0,ij,e_i^2=1`$. In $`V_1`$ we can find a basis $`v_1,w_1,\mathrm{},v_n,w_n`$ such that $`f(v_i,w_j)=\delta _{ij}`$, $`f(v_i,v_j)=f(w_i,w_j)=0`$. Consider the Weyl algebra $`W_n=1,v_i,w_i,1in,[v_i,w_j]=\delta _{ij},[v_i,w_j]=[v_i,w_j]=0`$. According to , the universal associative enveloping algebra of $`F+V`$ is isomorphic to the (super)tensor product $`C(V,f)_FW_n`$. We will utilize this fact in the following Lemma.
###### Lemma 2.4
There are no subsuperalgebras $``$ isomorphic to $`J(V,f)`$, where $`V=V_0+V_1`$, $`V_1\{0\}`$ in a finite-dimensional Jordan superalgebra $`๐^{(+)}`$, where $`๐`$ is an associative superalgebra.
Proof. We assume the contrary, that is, there exists a subsuperalgebra $``$ of the type $`J(V,f)`$ in $`๐^{(+)}`$. For $``$, we consider the universal associative enveloping superalgebra $`U()`$. According to the above fact, $`U()=C(V_0,f)_FW_n`$ where $`C(V_0,f)`$ is a Clifford algebra for $`V_0`$, $`f`$ is a bilinear form on $`V_0`$, $`W_n`$ is a Weyl algebra, $`n=\frac{1}{2}dimV_1`$. Let $`\phi `$ denote the identity embedding of $``$ in $`๐`$. As a direct consequence of the definition of universal enveloping algebra, $`\phi `$ can be uniquely extended to a homomorphism $`\overline{\phi }:U()๐`$. Note that $`\overline{\phi }(x)=\phi (x)=x`$ where $`xV_1`$. In other words, $`\overline{\phi }(V_1)0`$. However, since $`V_1`$ generates $`W_n`$, $`\overline{\phi }(W_n)0`$. It follows that $`\overline{\phi }(W_n)W_n`$. Therefore, $`๐`$ has an infinite-dimensional subsuperalgebra. This contradicts our assumptions.
In the next Lemma we will prove that no simple decompositions in which one of the components has either the type $`K_3`$ or $`D_t`$ are possible.
###### Lemma 2.5
Any superalgebra $`๐ฅ`$ of the type $`M_{n,m}(F)^{(+)}`$, $`n,m>0`$ cannot be represented as the sum of two proper simple subsuperalgebras one of which has the type $`K_3`$ or $`D_t`$.
Proof. First of all, we note that if $`\text{char}F=0`$, then $`K_3`$ has no non-trivial finite-dimensional associative specializations (see Remark 1). Therefore, we can directly pass to the second case when one of the subsuperalgebras is isomorphic to $`D_t`$. Next we suppose that $`๐ฅ`$ is given in the canonical form which is the set of all matrices of order $`(n+m)`$ with respect to the natural $`Z_2`$-gradation. Let $`V=V_0+V_1`$ denote the $`Z_2`$-graded module corresponding to the natural representation of $`M_{n,m}(F)^{(+)}`$, $`dimV_0=n`$, $`dimV_1=m`$. Any decomposition of $`M_{n,m}(F)^{(+)}`$ induces that of $`(M_{n,m}(F)^{(+)})_0=H(_n)H(_m)`$ given by
$$H(_n)H(_m)=(Fe_1Fe_2)+_0,$$
where $`e_1`$, $`e_2`$ are pairwise orthogonal idempotents in $`๐D_t`$. Next we define a pair of homomorphisms denoted as $`\pi _1`$, $`\pi _2`$ which are the projections on the ideals $`H(_n)`$ and $`H(_m)`$, respectively. Then the above decomposition can be rewritten in the following way:
$$H(_n)=\pi _1(Fe_1Fe_2)+\pi _1(_0),$$
$$H(_m)=\pi _2(Fe_1Fe_2)+\pi _2(_0).$$
Next we estimate the dimension of $`\pi _1(_0)`$. We know that $`\pi _1(_0)`$ is either a simple or a non-simple semisimple subalgebra. Therefore, in the first case we have $`dim\pi _1(_0)n^22n+1`$, and in the second case $`dim\pi _1(_0)n^22n+2`$. It follows that $`dimH(_n)2+dim\pi _1(_0)`$, $`n^2n^22n+4`$, $`n2`$. Using the same arguments as above we can prove that $`m2`$. As a result, we have four possible cases $`๐ฅ=M_{1,1}(F)^{(+)}`$, $`M_{1,2}(F)^{(+)}`$, $`M_{2,1}(F)^{(+)}`$ or $`M_{2,2}(F)^{(+)}`$. Notice that the first case can be immediately excluded because $`dimM_{1,1}(F)^{(+)}=4`$. Since $`M_{1,2}(F)^{(+)}M_{2,1}(F)^{(+)}`$ the second and forth are the only cases of interest to us. Let $`๐ฅ=M_{1,2}(F)^{(+)}`$. By the dimension argument, $`dim5`$, and, moreover, $`\text{rk}3`$. Clearly, there are only three appropriate choices for $`:`$ $`osp(2,1)`$, $`P(2)`$ or $`Q(2)`$. However, for all cases, $`U()`$ is isomorphic to $`M_{2,2}(F)`$, that is, cannot be a subsuperalgebra of $`M_{1,2}(F)`$.
Finally, let $`๐ฅ=M_{2,2}(F)^{(+)}`$. Again by the dimension argument, $`dim12`$ and $`\text{rk}_04.`$ Considering all possible cases we come to the conclusion that there are no appropriate subsuperalgebras in $`๐ฅ`$. This proves our Lemma.
###### Lemma 2.6
Any superalgebra of the type $`M_{n,m}(F)^{(+)}`$, $`n,m>0,`$ cannot be represented as the sum of two proper non-trivial simple subsuperalgebras one of which has either the types $`M_{1,1}(F)^{(+)}`$, $`osp(1,1)`$ or $`P(2)`$.
Proof. Assume that $`M_{n,m}(F)^{(+)}=๐+`$ where $`๐`$ and $``$ satisfy all the above conditions. Let $`๐`$ have one of types $`M_{1,1}(F)^{(+)}`$, $`osp(1,1)`$ or $`P(2)`$. The even part of the above decomposition can be rewritten as follows:
$$H(_n)=\pi _1(๐_0)+\pi _1(_0),$$
$$H(_m)=\pi _2(๐_0)+\pi _2(_0).$$
Let $`๐`$ be isomorphic to either $`M_{1,1}(F)^{(+)}`$ or $`osp(1,1)`$. Then $`dim\pi _1(๐_0)=2`$. Acting in the same manner as in Lemma 2.5, we obtain the following possibilities: $`n=m=1`$, $`n=1`$, $`m=2`$ ($`m=2`$, $`n=1`$), $`n=m=2`$. Obviously, there are no possible simple decompositions in the first case due to the low dimension of $`M_{1,1}(F)^{(+)}`$. In the second and third cases we have the following restrictions on the dimension and the rank of $`_0`$: $`dim_05`$, $`\text{rk}_03`$; $`dim_012`$, $`\text{rk}_04`$. Considering all cases one after another we conclude that there is no suitable choice for $`_0`$. Therefore, $`M_{n,m}(F)^{(+)}๐+`$.
In the last case when $`๐P(2)`$ there are the following restrictions on indices: $`n3`$, $`m3`$. In other words, $`n=1`$, $`m=2`$; $`n=m=2`$; $`n=1`$, $`m=3`$; $`n=2`$, $`m=3`$; $`n=m=3`$. By the dimension and rank arguments there is no such $`_0`$. The Lemma is proved.
Next taking into account all previous Lemmas we list simple decompositions that might exist in $`M_{n,m}(F)^{(+)}`$. Let $`๐`$ and $``$ stand for the simple non-trivial Jordan subsuperalgebras of $`M_{n,m}(F)^{(+)}`$.
$$\left|\begin{array}{ccc}& & \\ & ๐& \\ & & \\ 1& M_{k,l}(F)^{(+)}& M_{p,q}(F)^{(+)}\\ 2& M_{k,l}(F)^{(+)}& P(q)\\ 3& M_{k,l}(F)^{(+)}& Q(p)\\ 4& P(k)& Q(l)\\ 5& P(k)& P(l)\\ 6& Q(k)& Q(l)\\ 7& osp(k,l)& M_{p,q}(F)^{(+)}\\ 8& osp(k,l)& Q(p)\\ 9& osp(k,l)& P(q)\\ 10& osp(k,l)& osp(p,q)\end{array}\right|$$
Considering associative subalgebras $`S(๐)`$ and $`S()`$ generated by $`๐`$ and $``$, respectively, we obtain a new decomposition of the form $`M_{n+m}(F)=S(๐)+S()`$ where $`S(๐)`$ and $`S()`$ are associative subalgebras of $`M_{n+m}(F)`$. Note that $`S(๐)`$ is a homomorphic image of $`U(๐)`$. As a direct consequence of Theorem 2.2, $`U(๐)`$ is either an associative simple algebra or a direct sum of two or more simple pairwise isomorphic associative algebras.
###### Lemma 2.7
Let $`๐`$ be a proper non-trivial simple subsuperalgebra in $`M_{n,m}(F)^{(+)}`$ where $`n,m>0`$. Then $`S(๐)`$ coincides with $`M_{n+m}(F)`$ if and only if one of the following conditions hold
(1) Either $`๐osp(p,q)`$, $`p+2q=n+m`$, or
(2) $`๐P(n)`$ for the case when $`n=m`$.
Proof. First, we note that the converse of this Lemma is obvious (see Theorem 2.2). To prove that one of the above conditions holds in the case when $`S(๐)=M_{n,m}(F)`$, then we first show that $`๐`$ cannot be of type $`M_{k,l}(F)^{(+)}`$ or $`Q(p)`$. If $`๐`$ has the type $`M_{k,l}(F)^{(+)}`$, then $`k+l<n+m`$. By Theorem 2.2, $`S(๐)`$ is either a proper simple subalgebra of the type $`M_{k+l}(F)`$ or a non-simple semisimple subalgebra of the type $`M_{k+l}(F)M_{k+l}(F)`$. In both cases, $`S(๐)M_{n+m}(F)`$.
If $`๐Q(k)`$, then its associative enveloping algebra is a non-simple semisimple subalgebra which is the direct sum of two or more simple ideals of the type $`M_k(F)`$. Therefore, $`S(๐)M_{n+m}(F)`$.
For the rest cases, $`๐`$ can either have the type $`osp(p,q)`$ or $`P(k)`$. If $`๐osp(p,q)`$, then $`S(๐)M_{p+2q}(F)`$. Hence $`S(๐)=M_{n+m}(F)`$ if and only if $`p+2q=n+m`$. This yields (1).
Next we continue our proof by assuming that $`nm`$, say, $`n<m`$. We let $`๐`$ have the type $`P(k)`$. Then its even component $`๐_0`$, which is isomorphic to $`H(_k)`$, is a proper subalgebra in $`M_{n,m}(F)_0^{(+)}=I_1I_2`$, $`I_1H(_n)`$, $`I_2H(_m)`$. As previously, let $`\pi _1`$ and $`\pi _2`$ denote the projections on $`I_1`$ and $`I_2`$, respectively.
Suppose that $`S(๐)=M_{n+m}(F)`$. Since $`๐P(k)`$, then $`S(๐)M_{2k}(F)`$. This implies $`2k=n+m`$, $`k=\frac{n+m}{2}`$. In particular, $`k>n`$. Hence $`\pi _1(๐_0)=\{0\}`$ and $`\pi _2(๐_0)H(_k)`$. It follows that $`๐_0I_2`$. Thus the identity $`e`$ of $`๐`$ is an element of $`I_2`$. For any $`x๐_1`$, $`xe+ex=2x`$ where the multiplication is associative. Multiplying both sides of this equation by $`e`$, we obtain the following $`exe+ex=2ex`$. Since $`exe=0`$, we have $`ex=2ex`$. Similarly, $`xe=0`$, that is, $`x=0`$, for any $`x๐_1`$, a contradiction.
In conclusion, it remains to consider the case when $`n=m`$ and $`๐P(n)`$. However, it is obvious that $`S(๐)M_{2k}(F)`$ and $`S(๐)=M_{2n}(F)`$ if and only if $`k=n`$. This completes our proof.
###### Lemma 2.8
Let $`M_{n,m}(F)^{(+)}=๐+`$, $`n,m>0`$. Then one of the subsuperalgebras in the given decomposition has either the type $`osp(p,q)`$ where $`p+2q=n+m`$ or $`P(n)`$ (only if $`n=m`$).
Proof. Let us assume the contrary, that is, neither $`๐`$ nor $``$ is a subsuperalgebra of any of the above types. Then, by Lemma 2.7, $`S(๐)`$ and $`S()`$ are proper associative subalgebras in $`M_{n+m}(F)`$. Theorem 2.2 states that both $`S(๐)`$ and $`S()`$ are either simple associative algebras or non-simple semisimple associative algebras decomposable into the sum of two or more pairwise isomorphic simple algebras. Therefore, $`dimS(๐)k^2(\frac{n+m}{k})=(n+m)k`$ where $`k^2`$ is a dimension of a simple ideal, $`k>1`$. If one of the subsuperalgebras in the decomposition of $`M_{n+m}(F)`$ has a non-zero annihilator then by Proposition in no such decomposition exists. Hence the identity of $`M_{n+m}(F)`$ is contained in the intersection of $`S(๐)`$ and $`S()`$. On the other hand, the dimension of $`S(๐)`$ as well as $`S()`$ is strictly greater then $`\frac{(n+m)^2}{2}`$.
Thus, by the dimension argument, the sum of $`S(๐)`$ and $`S()`$ is a proper vector subspace of $`M_{n+m}(F)`$. Therefore, $`M_{n+m}(F)S(๐)+S()`$. This implies that our hypothesis was wrong.
###### Lemma 2.9
Let $`M_{n,m}(F)^{(+)}=๐+`$, $`n,m>0`$. Then, in the case when $`m`$ is even, and $`n`$ is odd, $`๐osp(n,\frac{m}{2})`$ and $`M_{k,l}(F)^{(+)}`$ where either $`k=n1,n`$ or $`l=m`$. On the contrary, if $`m`$ is odd, and $`n`$ is even, then $`๐osp(m,\frac{n}{2})`$ and $`M_{k,l}(F)^{(+)}`$ where either $`k=m1,m`$ or $`l=n`$.
Proof.
Since the proof remains the same for both cases, we consider only the first case. First, let $`nm`$. In view of Lemma 2.8, one of the subsuperalgebras in $`M_{n,m}(F)^{(+)}=๐+`$, for example $`๐`$, is isomorphic to $`osp(p,q)`$ where
$$p+2q=n+m.$$
$`(1)`$
The decomposition of $`M_{n,m}(F)^{(+)}`$ given above induces the following representation of the even component $`M_{n,m}(F)_0^{(+)}=๐_0+_0`$ where $`M_{n,m}(F)_0^{(+)}=H(_n)H(_m)`$, $`๐_0H(F_p)H(๐ฌ_q)`$. If for some $`i`$, $`\pi _i(๐_0)H(F_p)H(๐ฌ_q)`$, then either $`p+2qn`$ or $`p+2qm`$. However these inequalities conflict with condition (1). Hence either $`\pi _1(๐_0)H(F_p)`$, $`\pi _2(๐_0)H(๐ฌ_q)`$ or $`\pi _1(๐_0)H(๐ฌ_q)`$, $`\pi _2(๐_0)H(F_p)`$. If the first possibility holds true, then
1. $`H(_n)=\pi _1(๐_0)+\pi _1(_0)`$, $`\pi _1(๐_0)H(F_p)`$, $`pn`$, $`\pi _1(_0)0`$.
2. $`H(_m)=\pi _2(๐_0)+\pi _2(_0)`$, $`\pi _2(๐_0)H(๐ฌ_q)`$, $`q\frac{m}{2}`$, $`\pi _2(_0)0`$.
Since $`p+2q=n+m`$, it follows that $`p=n`$ and $`q=\frac{m}{2}`$. Clearly, $`๐`$ has the type $`osp(n,\frac{m}{2})`$. If the second possibility holds true, then acting in the same manner, we can show that $`p=m`$, $`q=\frac{n}{2}`$. However, we assumed that $`n`$ is odd. Hence it remains to prove that $`M_{k,l}(F)^{(+)}`$ where a pair of indices $`k,l`$ satisfies the conditions given in the Lemma. To prove this, we consider all possible types for $``$ in a step-by-step manner.
If $`๐osp(n,\frac{m}{2})`$, $`P(k)`$, then the decomposition induces the following representation of the odd part: $`M_{n,m}(F)_1^{(+)}=๐_1+_1`$ where $`dim๐_1=nm`$, $`dim_1=k^2`$, that is, $`2nmnm+k^2`$, $`nmk^2`$. Conversely, $`kn`$, $`km`$ since both projections $`\pi _1(_0)`$, $`\pi _2(_0)`$ are non-zero. Moreover, one of the inequalities should be strict since $`nm`$. Therefore, $`k^2<nm`$, which is a contradiction.
If $`๐osp(n,\frac{m}{2})`$, $`Q(k)`$, then, acting in the same manner as in the previous case, we can prove that $`M_{n,m}(F)^{(+)}๐+`$.
If $`๐osp(n,\frac{m}{2})`$, $`osp(p,q)`$, then, by the dimension argument, we have the following inequality $`(n+m)^22(\frac{n(n+1)}{2}+\frac{m(m1)}{2}+nm)`$. Simplifying the last inequality, we obtain that $`mn`$. Clearly, the opposite inequality $`mn`$ also holds true. Therefore, $`m=n`$ which contradicts our hypothesis. Overall, it remains to consider the case when $`๐osp(n,\frac{m}{2})`$, $`M_{k,l}(F)^{(+)}`$.
Again the decomposition of $`M_{n,m}(F)^{(+)}`$ induces that of $`M_{n,m}(F)_0^{(+)}`$ as follows: $`M_{n,m}(F)_0^{(+)}=๐_0+_0`$. Moreover, $`M_{n,m}(F)_0^{(+)}=H(_n)H(_m)`$, $`๐_0H(F_n)H(๐ฌ_{\frac{m}{2}})`$, $`_0H(_k)H(_l)`$. If both $`\pi _1(_0)`$ and $`\pi _2(_0)`$ are non-simple semisimple, that is, $`\pi _1(_0)_0`$ and $`\pi _2(_0)_0`$, then we have the following restrictions: $`k+ln`$ and $`k+lm`$. Since $`nm`$, we can assume without any loss of generality that $`n<m`$. Hence the dimension of $`\pi _i(_0)`$, $`i=1,2`$, is less than $`n^22n+2`$. It follows from $`\pi _1(_0)_0`$ that $`dim_0n^22n+2`$.
As a result, $`dimM_{n,m}(F)_0^{(+)}=n^2+m^2\frac{n(n+1)}{2}+\frac{m(m1)}{2}+n^22n+2`$, $`\frac{m(m+1)}{2}\frac{n(n+1)}{2}+22n`$, which is wrong. Therefore, we have only two possibilities: either $`\pi _1(_0)`$ or $`\pi _2(_0)`$ is a simple algebra. According to , for the first case, $`k=n1,n`$ and, for the second, $`l=m`$. Thus the Lemma is proved for the case when $`nm`$.
To complete our proof we consider the case when $`n=m`$. First, we assume that neither $`๐`$ nor $``$ has the type $`P(n)`$. By the previous Lemma one of the subsuperalgebras, for example $`๐`$, is isomorphic to $`osp(p,q)`$, $`p+2q=2n`$, that is, $`p=n`$, $`q=\frac{n}{2}`$. Then
$$H(_n)=\pi _1(๐_0)+\pi _1(_0),\pi _1(๐_0)H(F_n),\pi _1(_0)0$$
$$H(_n)=\pi _2(๐_0)+\pi _2(_0),\pi _2(๐_0)H(๐ฌ_{\frac{n}{2}}),\pi _2(_0)0.$$
For some $`i`$, let $`\pi _i(_0)`$ be a non-simple semisimple subalgebra, then
$$\pi _i(_0)\{\begin{array}{c}H(_k)H(_l),k+ln\text{or}\\ H(F_k)H(๐ฌ_l),k+2ln\end{array}$$
Therefore, $`dim\pi _i(_0)n^22n+2.`$ However $`dimM_{n,n}(F)_0^{(+)}dim๐_0+dim_0`$, $`2n^2n^22n+2+\frac{n(n+1)}{2}+\frac{n(n1)}{2}=2n^22n+2`$, that is, $`n1`$. As mentioned above, there are no simple decompositions in $`M_{1,1}(F)^{(+)}`$. Hence both $`\pi _1(_0)`$ and $`\pi _2(_0)`$ are simple. It follows that $`\pi _1(_0)H(_{n1})`$, $`\pi _2(_0)H(_n)`$, that is, $`_0M_{n1,n}(F)`$.
Next we let $`๐`$ be of the type $`P(n)`$. Then $`\{\begin{array}{c}P(k)\\ Q(k)\\ osp(k,l)\\ M_{k,l}(F)\end{array},`$ for some integers $`k`$ and $`l`$.
1. $`P(k)`$, hence $`dim=2k^2`$, $`kn`$. For $`dimM_{n,n}(F)^{(+)}dim๐+dim`$, it is clear that $`k=n`$, and the sum in the decomposition is direct. However since both subsuperalgebras have the type $`P(n)`$, they contain the identity of $`M_{n,n}(F)^{(+)}`$, a contradiction.
2. $`Q(k)`$. In this case the proof is the same as in Case 1.
3. $`osp(k,l)`$. Clearly, for some $`i`$, $`\pi _i(_0)`$ is non-simple semisimple. Therefore, $`k+2ln`$. In particular, $`\pi _i(_0)_0`$. Hence, $`dim_0n^22n+2`$. Thus $`dimM_{n,n}(F)_0^{(+)}=2n^22n^22n+2`$, a contradiction. As a result, $`\pi _i(_0)`$, $`i=1,2`$, is a simple subalgebra, that is, $`kn`$, $`l\frac{n}{2}`$. Then $`dim2n^2`$. By the dimension argument, $`k=n`$, $`l=\frac{n}{2}`$ and the sum in the given decomposition is direct. However, this contradicts the fact that both subsuperalgebras in the given decomposition contain the identity of $`M_{n,n}(F)^{(+)}`$.
4. $`M_{k,l}(F)^{(+)}`$, $`k+l<2n`$. The even part of $`M_{n,n}(F)^{(+)}`$, that is, $`M_{n,n}(F)_0^{(+)}`$ equals to the sum of two orthogonal ideals denoted as $`I_1`$ and $`I_2`$, both ideals isomorphic to $`H(_n)`$. By the dimension argument, $`dimM_{n,n}(F)^{(+)}2n^2+(k+l)^2`$, $`4n^22n^2+(k+l)^2`$, $`k+l\sqrt{2}n`$. In particular, $`\pi _1(_0)`$, $`\pi _2(_0)`$ are simple. Therefore, acting by a appropriate automorphism of $`M_{n,n}(F)^{(+)}`$, $`_0`$ can be reduced to the block-diagonal form. Moreover, $`I_1`$ and $`I_2`$ contain all simple ideals isomorphic to $`H(_k)`$ and $`H(_l)`$, respectively.
Suppose that the identity of $`M_{n,n}(F)^{(+)}`$ is an element of $``$. This implies that $`kk_1=ll_1=n`$ where $`k_1`$ and $`l_1`$ are the numbers of blocks which have types $`H(_k)`$ and $`H(_l)`$, respectively. In view of the inequality $`k+l\sqrt{2}n`$ this result implies that either $`k_1=2`$, $`l_1=1`$ or $`k_1=1`$, $`l_1=2`$, that is, $`M_{\frac{n}{2},n}(F)`$ up to the order of indices. By Theorem 2.2, $`S()`$ is isomorphic to $`M_{\frac{3n}{2}}(F)`$ or $`M_{\frac{3n}{2}}(F)M_{\frac{3n}{2}}(F)`$. Obviously, $`S()`$ cannot be non-simple semisimple of the indicated type because its rank is greater than $`2n`$. However, by Lemma 1.1, the first case is also impossible because $`\frac{3n}{2}>n`$. Hence the identity of $`M_{n,n}(F)`$ is not an element of $``$. In other words, $``$ as well as $`_0`$ has a non-zero annihilator.
Acting by appropriate automorphism of $`M_{n,n}(F)`$ we can reduce $``$ to the following form:
$$\left\{\left(\begin{array}{ccccccc}0& 0& \mathrm{}& 0& 0& \mathrm{}& 0\\ & & & & & & \\ 0& & & & & & \\ \mathrm{}& & T_1& & & T_2& \\ 0& & & & & & \\ & & & & & & \\ 0& & & & & & \\ \mathrm{}& & T_3& & & T_4& \\ 0& & & & & & \end{array}\right)\right\}$$
where $`T_1`$, $`T_2`$, $`T_3`$ and $`T_4`$ are matrices of orders $`(n1)\times (n1)`$, $`(n1)\times m`$, $`m\times (n1)`$ and $`m\times m`$, respectively.
This implies that $`๐_0`$ takes the form:
$$\left\{\left(\begin{array}{cc}X& 0\\ 0& C^1X^tC\end{array}\right)\right\},$$
for some $`C`$, $`\text{det}C0`$. Then using the automorphism $`\phi (Y)=C^1YC^{}`$ where
$$C^{}=\left(\begin{array}{cc}I& 0\\ 0& C^1\end{array}\right),$$
$`๐`$ can be reduced to the form where
$$๐_0=\left\{\left(\begin{array}{cc}X& 0\\ 0& X^t\end{array}\right)\right\}$$
while $``$ remains the same. Obviously, this decomposition is not possible. The Lemma is proved.
###### Example 1
A Jordan superalgebra of the type $`M_{n,m}(F)^{(+)}`$ where $`m`$ is even can be represented as the sum of two proper simple subsuperalgebras $`๐`$ and $``$ which have types $`osp(n,\frac{m}{2})`$ and $`M_{n1,m}(F)^{(+)}`$, respectively.
Proof. To prove, we consider the first subsuperalgebra in the standard realization:
$$\left\{\left(\begin{array}{cc}A& C\\ & \\ S^1C^t& B\end{array}\right)\right\}$$
where $`A`$ is a symmetric matrix of order $`n`$, $`B`$ is a symplectic matrix of order $`m`$, $`C`$ is any matrix of order $`n\times m`$, $`S=\left(\begin{array}{cc}0& I\\ I& 0\end{array}\right)`$ where $`I`$ is identity matrix of order $`\frac{m}{2}`$. The second subalgebra can be viewed in the following form:
$$\left\{\left(\begin{array}{ccccccc}& & & & & & \\ & A& & & B& & \\ & & & & & & \\ & & & & & & \\ 0& \mathrm{}& 0& 0& \mathrm{}& 0& \\ & & & & & & \\ & & & & & & \\ & C& & & D& & \end{array}\right)\right\}$$
where $`A`$ and $`C`$ of orders $`(n1)\times n`$ and $`m\times n`$, respectively, have the last two columns equal, $`B`$ and $`D`$ are any matrices of orders $`(n1)\times n`$, $`m\times m`$, respectively. By straightforward calculations $`dim(๐_1+_1)=dim๐_1+dim_1dim(๐_1_1)=mn+2m(n1)m(n2)=2mn`$. This proves our Lemma.
###### Example 2
In the case when $`n`$ is even, a Jordan algebra of the type $`M_{n,m}(F)^{(+)}`$ can also be decomposed into the sum of $`๐`$ and $``$ where $`๐osp(m,\frac{n}{2})`$ and $`M_{m1,n}(F)^{(+)}`$. This decomposition can be constructed in the same manner as in Example 1.
###### Proposition 1
Example 1 and 2 are the only possible decompositions of $`M_{n,m}(F)^{(+)}`$, $`n,m>0`$ into the sum of two proper simple non-trivial subsuperalgebras for appropriate values of $`n`$, $`m`$.
Proof. As usual, we assume the contrary, that is, there exists some other simple decomposition of $`M_{n,m}(F)^{(+)}`$ different from one in Example 1. By Lemma 2.9, this decomposition takes the following form:
1. If $`m`$ is even, then $`M_{n,m}(F)^{(+)}=๐+`$, $`๐osp(n,\frac{m}{2})`$, $`M_{l,k}(F)^{(+)}`$ where either $`l=n1,n`$ or $`k=m`$.
2. If $`n`$ is even, then $`M_{n,m}(F)^{(+)}=๐+`$, $`๐osp(m,\frac{n}{2})`$, $`M_{k,l}(F)^{(+)}`$ where either $`k=m1,m`$ or $`l=n`$.
Then $`M_{n,m}(F)_1=๐_1+_1`$. It follows that $`dimM_{n,m}(F)_1dim๐_1+dim_1`$, that is, $`2nmnm+2lk`$, $`nm2lk`$. Hence, for even $`m`$, $`l\frac{n}{2}`$, in the case $`k=m`$, and $`k\frac{m}{2}`$, in the case $`l=n1`$ or $`n`$. Likewise, if $`n`$ is even, then $`k\frac{m}{2}`$, in the case $`l=n`$, and $`l\frac{n}{2}`$, in the case $`k=m1`$ or $`m`$. For definitness, we consider the case when $`m`$ is even, and $`l=n1`$ because the proof remains the same for all other cases.
Let $`V=V_0+V_1`$ denote a $`Z_2`$-graded vector space where $`dimV_0=n`$ and $`dimV_1=m`$. By its definition, $`M_{n,m}(F)^{(+)}`$ coincides with the set of all linear transformations acting in $`V`$. Then let $`\rho `$ stand for the natural representation of $`=_0+_1`$ in $`V`$. It follows from the definition of this action that $`\rho (_0)(V_0)V_0`$, $`\rho (_0)(V_1)V_1`$, $`\rho (_1)(V_0)V_1`$, $`\rho (_1)(V_1)V_0`$. Since $``$ is a non-simple semisimple Jordan algebra it acts completely reducibly in $`V`$. Next we describe this action in more details. For this, we identify $`V`$ with a $`Z_2`$-graded vector space of the form $`W=v_0(V_0^{}F^{r+1})V_1^{}`$, $`r1`$ where $`v_0`$ is a vector in $`V_0`$ annihilated by $`_0`$, $`V_0^{}`$ is an invariant complementary subspace of $`v_0`$, $`\rho (_0)|_{V_0^{}}H(_{n1})`$, $`V_1^{}`$ is an invariant subspace of $`V_1`$ such that $`_0`$, $`\rho (_0)|_{V_1^{}}H(_k)`$. Moreover, $`W_0=v_0V_0^{}e_0`$, $`W_1=V_0^{}e_1,\mathrm{},e_rV_1^{}`$ where $`e_0,e_1,\mathrm{},e_r`$ is a basis for $`F^{r+1}`$. Then, $`\rho (_0)=\rho (_0)|_{v_0}\rho (_0)|_{V_0^{}}Id_{r+1}\rho (_0)|_{V_1^{}}.`$ Note that $`\rho (_0)|_{v_0}=0`$. In other words, by choosing an appropriate basis in $`V_0`$ and $`V_1`$, $`\rho (_0)`$ can be written in a block-diagonal form in which the first block of order 1 is zero, the last block has order $`k`$, and the other blocks have order $`r+1`$. Next we consider the representation of the odd part $`_1`$. For this, we choose any $`a_0`$ such that
$$\rho (a)(V_0^{}F^{r+1})=0,\rho (a)(V_1^{})0.$$
$`(2)`$
All such elements form an ideal of $`_0`$ isomorphic to $`H(_k)`$. Then we choose any non-zero $`x`$ in $`_1`$. Let $`e`$ denote the identity of $``$, $`e_0`$. Then $`\rho (x)v_0=\rho (xe)v_0=\rho (\frac{xe+ex}{2})v_0=\frac{1}{2}(\rho (x)\rho (e)v_0+\rho (e)\rho (x)v_0)=\frac{1}{2}\rho (x)v_0`$, that is, $`\rho (x)v_0=0`$, for any $`x_1`$. Next we find the representation of $`ax_1`$. As mentioned above, $`\rho (ax)(v_0)=0`$. Besides, $`2\rho (ax)(V_0^{}e_0)=\rho (a)\rho (x)(V_0^{}e_0)+\rho (x)\rho (a)(V_0^{}e_0)V_1^{}`$, $`\rho (ax)(V_0^{}e_1,\mathrm{},e_r)=0`$, $`\rho (ax)(V_1^{})V_0^{}e_0`$. Clearly, we can find $`c_0`$ whose action is given by the following formulae:
$$\rho (c)(V_0^{}F^{r+1})0,\rho (c)(V_1^{})=0.$$
$`(3)`$
Now we need to determine
$$c(xa).$$
$`(4)`$
Since $`2\rho (c(ax))=\rho (c)\rho (ax)+\rho (ax)\rho (c),`$ we have the following: $`\rho (c(ax))(v_0)=0`$, $`\rho (c(ax))(V_0^{}e_1,\mathrm{},e_r)=0`$. Besides,
$$\rho (c(xa))(V_0^{}e_0)=\rho (c)\rho (x)\rho (a)(V_0^{}e_0)+\rho (x)\rho (c)\rho (a)(V_0^{}e_0)+$$
$$\rho (a)\rho (c)\rho (x)(V_0^{}e_0)+\rho (a)\rho (x)\rho (c)(V_0^{}e_0)=\rho (a)\rho (x)\rho (c)(V_0^{}e_0)V_1^{}.$$
$`(5)`$
Similarly,
$$\rho (c(xa))(V_1^{})=\rho (c)\rho (x)\rho (a)(V_1^{})+\rho (x)\rho (c)\rho (a)(V_1^{})+\rho (a)\rho (c)\rho (x)(V_1^{})+$$
$$\rho (a)\rho (x)\rho (c)(V_1^{})=\rho (c)\rho (x)\rho (a)(V_1^{})V_0^{}e_0.$$
$`(6)`$
Assume that $`\rho (x)(V_1^{})0`$, $`\rho (x)(V_0^{}e_0)0(\text{mod}V_0^{}e_1,\mathrm{},e_r)`$. Then $`\rho (c(xa))`$ has the following matrix form:
$$\left(\begin{array}{cccc}0& \mathrm{}& 0\mathrm{}& 0\\ & & & \\ \mathrm{}& 0& 0\mathrm{}& XY_1Z\\ & & & \\ 0& 0& & \\ \mathrm{}& \mathrm{}& & \\ 0& ZY_2X& & 0\end{array}\right),$$
where $`X`$ is an arbitrary square matrix of order $`k`$, $`Y_1`$ and $`Y_2`$ are some fixed non-zero matrices of order $`k\times (n1)`$ and $`(n1)\times k`$, respectively, $`Z`$ is any square matrix of order $`n1`$. Next we choose any $`y_1`$. We have seen that there exists an element $`y`$ of form (4) such that $`\rho (ya(xc))(V_1^{})=0`$ or $`\rho (ya(xc))(V_0^{}e_0)=0(\text{mod}V_0^{}e_1,\mathrm{},e_r)`$. Suppose that one of the above equations does not hold. Without any loss of generality we let $`\rho (y^{})(V_0^{}e_0)0(\text{mod}V_0^{}e_1,\mathrm{},e_r)`$, where $`y^{}=ya(xc)`$. Multiplying $`y^{}`$ by the elements of the form (2) and then (3) we obtain $`a^{}(y^{}c^{})_1`$, where $`a^{}`$ and $`c^{}`$ ran relevant sets, and $`\rho (a^{}(y^{}c^{}))(V_0^{}F^{r+1})=0`$, $`\rho (a^{}(y^{}c^{}))V_1^{}=\rho (a^{})\rho (y^{})\rho (c^{})V_1^{}V_0^{}e_0`$. Moreover, $`\rho (a^{}(y^{}c^{})):V_1^{}V_0^{}e_0`$ represents all linear transformations from $`k`$-dimensional vector space into $`(n1)`$-dimensional vector space. Besides, all such elements are linearly independent from all the elements (4). Therefore, we found $`2(n1)k`$ linearly independent elements of $`_1`$, ($`dim_1=2(n1)k`$). If there is at least one element $`\overline{y}_1`$ such that $`\rho (\overline{y})(V_0^{}e_0)0(\text{mod}V_1^{})`$ or $`\rho (\overline{y})(V_0^{}e_1,\mathrm{},e_r)0`$, then it will be also linearly independent with all above elements. Hence, by dimension arguments, there is no $`\overline{y}`$ satisfying the above conditions. Consequently, for all elements in $`๐_1`$, the following $`\rho (\overline{y})(V_0^{})=0(\text{mod}V_1^{})`$, $`\rho (\overline{y})(V_0^{}e_1,\mathrm{},e_r)=0`$, $`\rho (\overline{y})(V_1^{})V_0^{}e_0`$ hold true. If we fix a basis in $`V`$ such that in this basis the even part has the diagonal form:
$$\left(\begin{array}{cccc}X& & 0& \\ & & & \\ & X& \mathrm{}& 0\\ 0& \mathrm{}& \mathrm{}& \\ & 0& & Y\end{array}\right),$$
then the odd part becomes the following:
$$\left(\begin{array}{cccc}0& 0& \mathrm{}0& Z\\ & & & \\ 0& & & \\ \mathrm{}& & 0& \\ 0& & & \\ Z^{}& & & \end{array}\right)$$
$`(7)`$
where $`Z`$, $`Z^{}`$ are any matrices of order $`(n1)\times k`$ and $`k\times (n1)`$, respectively. Then it follows from $`_1_1_0`$ that $`_1=0`$, a contradiction.
We henceforth assume that the equations $`\rho (ya(xc))(V_1^{})=0`$ and $`\rho (ya(xc))(V_0^{}e_0)=0(\text{mod}V_0^{}e_1,\mathrm{},e_r)`$ hold true simultaneously. Then multiplying $`ya(xc)`$ by the elements (4), we obtain some elements of $`_0`$ which act on $`V_1^{}`$ and $`V_0^{}e_1,\mathrm{},e_r`$ non-invariantly. Hence, $`ya(xc)=0`$. Therefore, the odd component of $`๐_1`$ has form (7). As proved before, this is not possible.
Next we assume that $`\rho (x)(V_0^{}e_0)=0(\text{mod}V_0^{}e_1,\mathrm{},e_r)`$, for all $`x_1`$, and for at least one element $`x^{}_1`$, $`\rho (x^{})(V_1^{})0`$.
Acting in the same manner as before, we obtain $`a^{}(x^{}c^{})_1`$ which acts trivially on all subspaces except for $`V_1^{}`$, which it carries into $`V_0^{}e_0`$. Considering the difference between an arbitrary element $`y_1`$ and a corresponding element $`a^{\prime \prime }(x^{\prime \prime }c^{\prime \prime })`$, we can show that $`\rho (ya^{\prime \prime }(x^{\prime \prime }c^{\prime \prime }))(V_0^{}e_0)=0(\text{mod}V_0^{}e_1,\mathrm{},e_rV_1^{})`$, $`\rho (ya^{\prime \prime }(x^{\prime \prime }c^{\prime \prime }))(V_1^{})=0`$. Again multiplying $`a^{}(x^{}c^{})`$ and $`ya^{\prime \prime }(x^{\prime \prime }c^{\prime \prime })`$, we obtain some elements from $`_0`$ acting on $`V_1^{}`$ non-trivially. Then we conclude that $`_1`$ consists of all elements which act on $`V_0^{}e_0`$ trivially and carry the other subspaces into $`V_0^{}e_0`$. Hence $`_1_1=0`$, a contradiction.
Finally, if $`\rho (x)(V_1^{})=0`$, $`\rho (x)(V_0^{}e_0)=0(\text{mod}(V_0^{}F^{r+1}))`$, then it follows that $`_1_0=0`$, which is clearly a wrong statement. The Proposition is proved.
Based on all above Lemmas and Proposition 1, we conclude that Theorem 1 is true. In other words, $`M_{n,m}(F)^{(+)}`$ where $`n,m>0`$, $`m`$ is even, and $`n`$ is odd admits only one decomposition into the sum of two proper simple subsuperalgebras. If both $`n,m`$ are odd, then $`M_{n,m}(F)^{(+)}`$ cannot be represented as the sum of two proper simple subsuperalgebras. If both indices are even, then $`๐`$ admits two different types of decompositions of the following forms:
1. $`๐=_1+๐_1`$ where $`_1`$ and $`๐_1`$ have types $`osp(n,\frac{m}{2})`$ and $`M_{n1,m}(F)^{(+)}`$,
2. $`๐=_2+๐_2`$ where $`_2`$ and $`๐_2`$ have types $`osp(m,\frac{n}{2})`$ and $`M_{m1,n}(F)^{(+)}`$.
## 3 Decompositions of superalgebras of the type $`osp(n,m)`$
This section is dedicated to the study of simple decompositions of $`osp(n,m)`$. Actually, we will show that there are no such decompositions over algebraically closed field $`F`$ of zero characteristic. Our main purpose is to prove the following.
###### Theorem 3.1
Let $`๐ฅ`$ be a superalgebra of the type $`osp(n,m)`$ where $`n,m>0`$. Then $`๐ฅ`$ cannot be written as the sum of two proper nontrivial simple subsuperalgebras $`๐`$ and $``$.
The proof of this Theorem is based on the following Lemmas.
###### Lemma 3.2
Let $`๐ฅ`$ be a superalgebra of type $`osp(n,m)`$ where $`n,m>0`$, and $`๐`$, $``$ are two proper simple subsuperalgebras none of which has any of the types $`K_3`$ or $`D_t`$. Then $`๐ฅ`$ cannot be represented as the sum of $`๐`$ and $``$.
Proof. First we identify $`๐ฅ`$ with $`osp(n,m)`$ which can be considered in the canonical form. Next we assume the contrary, that is,
$$osp(n,m)=๐+,$$
$`(8)`$
The decomposition (8) generates the following decomposition of the associative enveloping algebra into the sum of three non-zero subspaces.
$$M_{n+2m}(F)=S(osp(n,m))=S(๐)+S()+S(๐)S(),$$
$`(9)`$
where $`S(๐)`$, $`S()`$ denote the associative enveloping algebras of $`๐`$, $``$, respectively. Let 1 denote the identity of $`osp(n,m)`$. Then we consider the following cases.
Case 1. Let $`1๐`$, $`1`$. This implies that there exist non-zero $`a_0`$ and $`b_0`$ in $`\text{Ann}(๐)`$ and $`\text{Ann}()`$, respectively. Then multiplying every term of (9) by $`a_0`$ on the left and $`b_0`$ on the right, the following equation $`a_0M_{n+2m}(F)b_0=0`$ takes place, which is clearly wrong.
Case 2. $`1๐`$, $`1`$. Six cases arise:
(a) $`๐M_{k,l}(F)^{(+)}`$, $`M_{p,q}(F)^{(+)}`$. The given decomposition induces the decomposition of the even part $`osp(n,m)_0=๐_0+_0`$ which in turn can be projected on the ideals of the even component. In particular, $`H(F_n)=\pi _1(๐_0)+\pi _1(_0)`$. By Theorem 1.3, both projections cannot be simultaneously simple. Therefore, at least one of the components is non-simple semisimple. For definiteness, let $`\pi _1(๐_0)H(_k)H(_l)`$. By Lemma 1.1, $`k+l\frac{n}{2}`$. Then, $`dim\pi _1(๐_0)=k^2+l^2(\frac{n}{2}1)^2+1=\frac{n^2}{4}n+2`$. If $`\pi _1(_0)H(_p)`$ or $`H(_q)`$, then, by Lemma 3.3, $`p\frac{n}{2}`$ or $`q\frac{n}{2}`$. If $`\pi _1(_0)H(_p)H(_q)`$, then $`dim\pi _1(_0)=p^2+q^2(\frac{n}{2}1)^2+1=\frac{n^2}{4}n+2`$. As a result, $`dimH(F_n)=\frac{n^2+n}{2}2(\frac{n^2}{4}n+2)`$, $`\frac{5n}{2}4`$, $`n1`$ or $`dimH(F_n)=\frac{n^2+n}{2}\frac{n^2}{4}n+2+\frac{n^2}{4}`$, $`\frac{3n}{2}2`$, $`n1`$.
There remains one case where $`n=1`$. The decomposition takes the following form: $`osp(1,m)=๐+`$ where $`๐M_{1,l}(F)^{(+)}`$, $`M_{p,q}(F)^{(+)}`$. Then $`H(๐ฌ_m)=\pi _2(๐_0)+\pi _2(_0)`$. If $`\pi _2(๐_0)FH(_l)`$ $`\pi _2(_0)H(_p)H(_q)`$, then $`1+lm`$, $`p+qm`$, $`dim\pi _2(๐_0)m^22m+2`$, $`dim\pi _2(_0)m^22m+2`$. As a result, $`2m^24m+4m(2m1)=2m^2m`$, $`43m`$, $`m1`$, that is, $`m=1`$. However, it is clear that both subsuperalgebras in $`osp(1,1)=๐+`$ are isomorphic to $`M(1,1)^{(+)}`$, and, by dimension argument, $`๐`$, $``$ coincide with $`M(1,1)^{(+)}`$.
If one of the projections is non-simple semisimple and the other is simple, then $`m(2m1)m^2+m^22m+2=2m^22m+2`$, $`m2`$. Therefore, $`osp(1,2)=๐+`$, $`๐M_{1,1}(F)^{(+)}`$, $`M_{1,2}(F)^{(+)}`$, $`H(๐ฌ_2)=\pi _2(๐_0)+\pi _2(_0)`$ where $`\pi _2(๐_0)`$ is simple, and $`\pi _2(_0)`$ is non-simple. By the dimension argument, the sum in the above decomposition is direct, that is, one of the subalgebras does not contain the identity, which is obviously wrong. If both projections are simple, then $`osp(1,m)=๐+`$, $`๐`$, $`M(1,m)^{(+)}`$. However, by the dimension argument, the latter does not hold. Otherwise, $`๐_1=_1=osp(1,m)_1`$ because their dimensions are equal. It follows that $`๐_0=๐_1๐_1=_1_1=_0`$, that is, $`๐==osp(1,m)`$.
(2)$`๐M_{k,l}(F)^{(+)}`$, $`P(q)`$ or $`Q(q)`$ ($`q>1`$). Therefore, $`H(F_n)=\pi _1(๐_0)+\pi _1(_0)`$ where $`\pi _1(๐_0)H(_k)H(_l)`$, $`\pi _1(_0)H(_q)`$. Again, by the same arguments as in the previous case, $`n1`$. Let $`n=1`$ then $`\pi _1(๐_0)F`$. In this case, $`k=1`$ (or $`l=1`$) and the following decomposition holds true: $`H(๐ฌ_m)=\pi _2(๐_0)+\pi _2(_0)`$ where $`\pi _2(๐_0)FH(_l)`$, $`\pi _2(_0)H(_q)`$ or $`\pi _2(๐_0)H(_l)`$, $`\pi _2(_0)H(_q)`$.
In the first case, we have proved that $`m=2`$, $`l=1`$. Hence $`osp(1,2)=๐+`$ where $`๐M_{1,1}(F)^{(+)}`$, $`P(2)`$, which induces the following: $`H(๐ฌ_2)=\pi _2(๐_0)+\pi _2(_0)`$ where $`๐_0FF`$, $`_0H(_2)`$. The sum in the last decomposition is direct, and both subalgebras contain 1, which is a contradiction. In the second case, $`1=q=m`$, that is, $`osp(1,m)=๐+`$, $`๐M(1,m)^{(+)}`$, $`P(m)`$, $`m>1`$.
Since $`dimosp(1,m)_1dim_1`$, then $`2mm^2`$, that is, $`m=2`$. If $`m=2`$, then $`osp(1,2)=๐+`$ where $`๐M_{1,2}(F)^{(+)}`$, $`P(2)`$ which induces the equality $`H(๐ฌ_2)=\pi _2(๐_0)+\pi _2(_0)`$ where $`๐_0H(_2)`$, $`_0H(_2)`$.
Notice that the identity of $`osp(1,2)`$ is an element of $`๐`$, that is, $`๐`$ has trivial two-sided annihilator. Consider an associative enveloping algebras of $`osp(1,2)`$ and $`๐`$ denoted as $`S(osp(1,2))`$ and $`S(๐)`$, respectively. It can be shown that $`S(๐)M_3(F)`$ is a subalgebra of $`S(osp(1,2))=M_5(F)`$ (see Theorem 2.2). By Lemma 1.1, $`S(๐)`$ contains no identity of $`M_5(F)`$, therefore, has a non-zero two-sided annihilator, and so does $`๐`$, a contradiction.
(c) $`๐`$, $``$ have types $`P(q)`$ or $`Q(p)`$. Then the decomposition leads to the decomposition of $`H(F_n)`$ into the sum of two proper subalgebras, which does not exist.
(d) $`๐osp(k,l)`$, $`M_{p,q}(F)^{(+)}`$. Since $`S(๐)M_{k+2l}(F)`$ contains the identity of the entire superalgebra, $`k+2l\frac{n+2m}{2}`$. Similarly, $`p+q\frac{n+2m}{2}`$.
Thus $`dimosp(n,m)dim๐+dim`$, that is, $`\frac{n^2+n}{2}+m(2m1)+2nm\frac{k^2+k}{2}+l(2l1)+2kl+\frac{(n+2m)^2}{4}`$. By straightforward calculations we obtain $`\frac{n^2}{4}+\frac{n}{2}+m^2+nm3m`$, which is true if and only if $`m=n=1`$. Obviously, $`osp(1,1)`$ has no simple decompositions.
(e) $`๐osp(k,l)`$, $`P(q)`$. Then, we have $`k+2l\frac{n+2m}{2}`$, $`2q\frac{n+2m}{2}`$. Therefore, $`dim=2q^22(\frac{n+2m}{4})^2`$. Again, by the dimension argument, this decomposition is not possible.
(f) $`๐osp(k,l)`$, $`osp(p,q)`$. Then $`k+2l\frac{n+2m}{2}`$, $`p+2q\frac{n+2m}{2}`$. Comparing $`dimosp(n,m)`$ with $`dim๐+dim`$ we have $`\frac{n^2}{2}+2nm+2m^24m`$, a contradiction.
Case 3 Let $`1๐`$, $`1`$. As mentioned above, the given decomposition induces the following decompositions of the ideals of the even component:
$$H(F_n)=\pi _1(๐_0)+\pi _1(_0),$$
$`(10)`$
$$H(๐ฌ_m)=\pi _2(๐_0)+\pi _2(_0).$$
$`(11)`$
If either $`\pi _1(๐_0)`$, $`\pi _2(_0)`$ or $`\pi _1(_0)`$, $`\pi _2(๐_0)`$ are non-simple semisimple, then $`dim๐_0=dim\pi _1(๐_0)<dimH(F_n)`$, $`dim_0=dim\pi _2(_0)<dimH(๐ฌ_m)`$. This implies that $`dim๐_0+dim_0<dim(H(F_n)H(๐ฌ_m))`$, which is wrong. Likewise we have a contradiction in the second case. Therefore, there is a simple algebra in each pair:($`\pi _1(๐_0)`$, $`\pi _2(_0)`$), ($`\pi _1(_0)`$, $`\pi _2(๐_0)`$). Since 1 is not an element of $``$, $``$ has a non-zero two-sided annihilator, and so does $`_0`$. It follows that one of $`\pi _1(_0)`$, $`\pi _2(_0)`$ has a non-zero two-sided annihilator. Let us assume the first possibility, that is, $`\pi _1(_0)`$ can be embedded in the simple subalgebra which also has a non-zero annihilator. Since $`H(F_n)`$ cannot be written as the sum of two simple subalgebras, $`\pi _1(๐_0)`$ should be non-simple semisimple. This implies that
$$\pi _1(๐_0)\{\begin{array}{c}H(F_k)H(๐ฌ_l),\text{ }\\ H(_k)H(_l)\end{array}$$
$`(12)`$
In other words, we represent $`H(F_n)`$ as the sum of a non-simple semisimple subalgebra of form (12) and a subalgebra which has a non-zero two-sided annihilator.
Let $`V`$ denote the $`n`$-column vector space. Then, there exists a non-zero vector $`vV`$ annihilated by the second subalgebra. By Lemma 1.2, $`dimH(F_n)v=n`$. It follows from (10) that $`dim\pi _1(๐_0)v=n`$.
If $`\pi _1(๐_0)H(F_k)H(๐ฌ_l)`$, then by some automorphism of $`F_n`$ it can be reduced to the following form:
$$\left(\begin{array}{cccccc}X& \mathrm{}& 0& 0& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& \mathrm{}& X& 0& \mathrm{}& 0\\ 0& \mathrm{}& 0& Y& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& \mathrm{}& 0& 0& \mathrm{}& Y\end{array}\right),$$
where $`X`$ is a symmetric matrix of order $`k`$, $`Y`$ is a symplectic matrix of order $`2l`$. Let $`v`$ be equal to $`(v_{11},\mathrm{},v_{1k_1},v_{21},\mathrm{},v_{2l_1})^t`$ where $`v_{i1}`$ is a vector of dimension $`k`$, $`i=1,\mathrm{},k_1`$, $`v_{2j}`$ is a vector of dimension $`2l`$, $`j=1,\mathrm{},l_1`$. Since $`\pi _1(๐_0)`$ contains 1, $`kk_1+2ll_1=n`$. Then, $`dim\{Xv_{1i}|XH(F_k)\}=k`$, $`dim\{Yv_{2j}|YH(๐ฌ_l)\}=2l1`$ (see Lemma 1.4). Therefore, $`dim\pi _1(๐_0)v=kk_1+(2l1)l_1<n`$, a contradiction. If $`\pi _1(๐_0)H(_k)H(_l)`$, then by some automorphism of $`F_n`$ it can be reduced to
$$\left(\begin{array}{ccccc}X& 0& \mathrm{}& 0& 0\\ 0& X^t& \mathrm{}& 0& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& \mathrm{}& Y& 0\\ 0& 0& \mathrm{}& 0& Y^t\end{array}\right).$$
$`(13)`$
Hence, $`\pi _2(_0)`$ has a non-trivial two-sided annihilator, that is, can be embedded in the simple algebra with a non-zero annihilator. Therefore, $`\pi _2(๐_0)`$ is non-simple semisimple because $`H(๐ฌ_m)`$ cannot be written as the sum of two simple subalgebras one of which has a non-zero two-sided annihilator (see ). As a result, we have the decomposition of the form: $`H(๐ฌ_m)=\pi _2(๐_0)+\pi _2(_0)`$ which in turn induces the following
$$F_{2m}=F_{2m1}+\pi _2(๐_0),$$
in which the first subalgebra clearly has a non-zero two-sided annihilator, and the second is non-simple semisimple. According to , such decomposition cannot exist. The Lemma is proved.
###### Lemma 3.3
A superalgebra $`๐ฅ`$ of the type $`osp(n,m)`$ where $`n,m>0`$ cannot be decomposed into the sum of two proper simple subsuperalgebras one of which has either the type $`K_3`$ or $`D_t`$.
Proof. First we identify $`๐ฅ`$ with $`osp(n,m)`$. Since Kaplanskyโs superalgebra $`K_3`$ has no finite-dimensional associative specializations, $`K_3`$ cannot be a subsuperalgebra of a superalgebra of the type $`osp(n,m)`$. Assume that $`osp(n,m)=๐+`$ where, for example, $`๐D_t`$. Then, the above decomposition induces the following:
$$H(F_n)=\pi _1(Fe_1Fe_2)+\pi _1(_0),$$
$$H(๐ฌ_m)=\pi _2(Fe_1Fe_2)+\pi _2(_0).$$
Let us note that $`dim\pi _1(_0)\frac{n(n1)}{2}`$ if it is a simple subalgebra and $`dim\pi _1(_0)\frac{n^23n+2}{2}+2`$ if it is a non-simple semisimple subalgebra. This implies that $`dimH(F_n)=\frac{n(n+1)}{2}4+\frac{n^2}{2}\frac{3n}{2}+1`$, $`2n5`$, $`n2`$. Similarly, $`dim\pi _2(_0)2m^25m+3`$ if $`\pi _2(_0)`$ is a simple subalgebra, and $`dim\pi _2(_0)2m^25m+4`$ if $`\pi _2(_0)`$ is a non-simple semisimple subalgebra. Thus $`dimH(๐ฌ_m)=2m^2m2+2m^25m+4`$, $`2m6`$, $`m\frac{3}{2}`$. Therefore, either $`๐ฅosp(1,1)`$ or $`๐ฅosp(2,1)`$. Since $`dimosp(1,1)=4`$, the first case is not possible. Let $`๐`$ be isomorphic to $`D_t`$. By the dimension argument, either $`M_{1,1}(F)`$ or $`osp(1,1)`$.
In turn, $`๐`$ acts completely reducibly in the 4-dimensional column vector space $`V`$ because $`dimV_0=dimV_1=2`$ (see Theorem 2.3). Moreover, $`V=W_1W_2`$, $`dimW_1=1`$, $`dimW_2=3`$, $`W_1`$, $`W_2`$ are invariant subspaces with respect to the action of $`๐`$. Besides, $`๐`$ acts in $`W_1`$ trivially and in $`W_2`$ irreducibly. It follows that, by some graded automorphism of $`osp(n,m)`$, $`๐`$ can be reduced to the following form:
$$๐=\left\{\left(\begin{array}{cccc}0& & \mathrm{}& 0\\ & & & \\ & & & \\ \mathrm{}& & A^{}& \\ 0& & & \end{array}\right)\right\}.$$
Since $`dim๐=4`$, $`๐osp(1,1)`$. Finally, we obtain a decomposition of a superalgebra of the type $`osp`$ into the sum of two subsuperalgebras of the same type. As proved before, such decomposition does not exist. The Lemma is proved.
## 4 Decompositions of superalgebras of types $`Q(n)`$ and $`P(n)`$
First of all we recall the canonical realizations of Jordan superalgebras of types $`P(n)`$ and $`Q(n)`$. A Jordan superalgebra of the type $`Q(n)`$ can be represented as the set of all matrices of order $`2n`$ which have the following form:
$$\left\{\left(\begin{array}{cc}A& B\\ B& A\end{array}\right)\right\}$$
where $`A`$ and $`B`$ are any square matrices of order $`n`$. A canonical realization of a Jordan algebra of the type $`P(n)`$ consists of all matrices of the form:
$$\left\{\left(\begin{array}{cc}A& B\\ C& A^t\end{array}\right)\right\}$$
where $`A`$ is any square matrix of order $`n`$, $`B`$ is a symmetric matrix of order $`n`$, $`C`$ is a skewsymmetric matrix of order $`n`$. Notice here that the even part of $`Q(n)`$ as well as $`P(n)`$ is isomorphic to a simple Jordan algebra of the type $`H(_n)`$. Besides, $`dimQ(n)=dimP(n)=2n^2`$. Later on only these two properties will be primarily used, hence, all Lemmas proved in this section are true for Jordan superalgebras of both types. For definiteness, we consider only Jordan subalgebras of type $`Q(n)`$. In several steps, we prove that no Jordan superalgebras of types $`P(n)`$ or $`Q(n)`$ can be represented as the sum of two proper simple subsuperalgebras.
###### Lemma 4.1
Let $`๐`$ of the type $`osp(p,q)`$ be a proper subsuperalgebra of $`Q(n)`$. Then $`dim๐\frac{n^2+n}{2}`$.
Proof. It follows from the Lemma conditions that $`๐_0H(F_p)H(๐ฌ_q)`$ is a proper subalgebra of $`Q(n)_0`$ which is isomorphic to $`H(_n)`$. Therefore, $`p+2qn`$, $`p,q>0`$. It is easy to see that the subalgebra takes on its maximum value when $`p+2q=n`$. Then $`dim๐=\frac{p^2+p}{2}+q(2q1)+2pq`$. This implies that $`dim๐\frac{n^2+n}{2}`$. The Lemma is proved.
###### Lemma 4.2
Let $`๐`$ of the type $`M_{k,l}(F)^{(+)}`$ where $`k,l>0`$ be a proper subsuperalgebras of $`Q(n)`$. Then $`dim๐n^2`$.
Proof. Since $`๐`$ is proper, $`k+ln`$, hence $`(k+l)^2n^2`$.
###### Lemma 4.3
A superalgebra $`๐ฅ`$ of either the type $`P(n)`$ or $`Q(n)`$, $`n>1`$, cannot be represented as the sum of two proper nontrivial subsuperalgebras one of which has either the type $`K_3`$ or $`D_t`$.
Proof. For definiteness, we assume that $`๐ฅ`$ has the type $`P(n)`$. Next, in order to simplify our notation we identify $`๐ฅ`$ with its canonical realization denoted as $`P(n)`$. By Remark 1 in Section 2 no superalgebra of the type $`K_3`$ can be a subsuperalgebra of $`P(n)`$. Therefore, let $`๐`$ be of type $`D_t`$. The given decomposition of $`P(n)`$ induces that of the form: $`P(n)_0=๐_0+_0`$ where $`๐_0=Fe_1Fe_2`$, $`e_1`$ and $`e_2`$ are pairwise orthogonal idempotents. Next we estimate the dimension of $`_0`$. If $`_0`$ is simple, then $`dim_0n^22n+1`$. If $`_0`$ is non-simple semisimple , then $`dim_0n^22n+2`$. As a result, $`dimP(n)_0=n^22+n^22n+2`$, $`n2`$. The only case which remains to prove is when $`n=2`$. By the dimension and rank arguments, either $`M_{1,1}(F)`$ or $`osp(1,1)`$. In both cases, $`_0`$ is isomorphic to $`Fe_1^{}Fe_2^{}`$ where $`e_1^{}`$ and $`e_2^{}`$ are pairwise orthogonal idempotents. Hence $`P(2)_0=H(_2)=Fe_1Fe_2+Fe_1^{}Fe_2^{}`$. Both subalgebras contain the identity. Thus, $`dimH(_2)=4>dim๐+dim`$. The Lemma is proved.
###### Lemma 4.4
Let $`๐ฅ`$ of the type $`P(n)`$ or $`Q(n)`$ be represented as the sum of two proper non-trivial subsuperalgebras $`๐`$ and $``$ whose even components are semisimple Jordan algebras and one of them has a non-zero two-sided annihilator. Then $`๐ฅ๐+`$.
Proof. Let $`๐ฅ=๐+`$, and $`๐`$ have a two-sided annihilator. Then $`๐ฅ_0=๐_0+_0`$ where $`๐_0`$, $`_0`$ are semisimple Jordan subalgebras, $`\text{Ann}๐_0\{0\}`$. Since $`๐ฅ_0H(_n)`$, $`๐ฅ_0`$ can be represented as the set of all matrices of order $`n`$, denoted as $`F_n^{(+)}`$, under the Jordan multiplication. Obviously, $`F_n=๐_0+_0`$ where $`๐_0`$ and $`_0`$ denote associative enveloping algebras for $`๐_0`$ and $`_0`$, respectively. This implies that $`F_n`$ can be written as the sum of two semisimple subalgebras $`๐_0`$ and $`_0`$ one of which has a non-zero two-sided annihilator. This contradicts Proposition 1 in . The Lemma is proved.
The following table summarizes all the information obtained above.
$$\left|\begin{array}{ccc}& & \\ & ๐& Maxdim\\ & & \\ 1& M_{k,l}(F)^{(+)}& n^2\\ 2& osp(p,q)& \frac{n^2+n}{2}\\ 3& Q(k)& 2(n1)^2\\ 4& P(k)& 2(n1)^2\end{array}\right|$$
In the second column we list all possible types which subsuperalgebras of $`P(n)`$ and $`Q(n)`$ can have. In the third column we point out the maximal dimension corresponding to each subsuperalgera.
###### Theorem 4.5
Let $`๐ฅ`$ have type either $`Q(n)`$ or $`P(n)`$, where $`n>1`$. Then $`๐ฅ`$ cannot be represented as the sum of two proper simple non-trivial subsuperalgebras $`๐`$ and $``$.
Proof. Since the case $`P(n)`$ is completely similar to the case $`Q(n)`$, we give proof only for $`Q(n)`$. We have the following cases.
Case 1. $`Q(n)=๐+`$, $`๐M_{k,l}(F)^{(+)}`$, $`M_{s,t}(F)^{(+)}`$. Besides, the dimensions of both subsuperalgebras are not greater than $`n^2`$. This implies that the sum in the above decomposition is direct. As a consequence of this fact, the given decomposition induces the decomposition of $`Q(n)_0`$ into the direct sum of two semisimple subalgebras $`๐_0`$ and $`_0`$ one of which does not contain the identity of the whole superalgebra or, equivalently, has a non-trivial two-sided annihilator. By Lemma 4.4, no such decomposition is possible.
Case 2. $`Q(n)=๐+`$, $`๐osp(p,q)`$, $`\{\begin{array}{c}osp(p^{},q^{})\\ M_{k,l}(F)^{(+)}\\ Q(k)\\ P(k)\end{array}.`$
Taking into account Lemma 4.1, we can conclude that the decomposition into the sum of two subsuperalgebras of the type $`osp`$ is not possible. Assume that the second decomposition holds true. According to the above estimates, the dimension of $`๐`$ and $``$ are not greater than $`\frac{n^2+n}{2}`$ and $`n^2`$, respectively. Hence, by the dimension argument, it is not possible. For the last two cases, the decomposition of the even part has the form: $`H(_n)=๐_0+_0`$, $`๐_0H(F_p)H(๐ฌ_q)`$, $`_0H(_k)`$. If $`1_0`$, then $`k\frac{n}{2}`$ and $`dim(๐+)\frac{n^2+n}{2}+\frac{n^2}{2}<2n^2`$. If $`1H(_k)`$, then we have a contradiction with Lemma 4.4.
Case 3. $`Q(n)=๐+`$, $`๐M_{k,l}(F)^{(+)}`$, $`\{\begin{array}{c}Q(m)\\ P(m)\end{array}.`$
This decomposition induces that of the even part: $`H(_n)=๐_0+_0`$, $`๐_0H(_k)H(_l)`$, $`_0H(_m)`$. If $`1_0`$, then again we have a contradiction with Lemma 4.4. If $`1_0`$, then $`k\frac{n}{2}`$, that is, $`dim\frac{n^2}{2}`$. However $`dim(๐+)n^2+\frac{n^2}{2}<2n^2`$, which is wrong.
Case 4. Let $`Q(n)=๐+`$, $`๐P(k)`$, $`Q(l)`$, $`k,l<n`$. As above, this decomposition induces the decomposition of the even part $`H(_n)`$ into the sum of two subalgebras of types $`H(_k)`$ and $`H(_l)`$. However it follows from the classification of simple decompositions of simple Jordan algebras that no such decomposition exists. The Theorem is proved.
## 5 Decompositions of superalgebras of types $`J(V,f)`$, $`K_3`$, $`D_t`$
###### Theorem 5.1
Let a superalgebra $`๐ฅ`$ have type $`J(V,f)`$, and $`๐`$, $``$ be simple non-trivial subsuperalgebras of $`๐ฅ`$. Then $`๐ฅ=๐+`$ implies that $`๐`$, $``$ are isomorphic to some superalgebras of non-singular symmetric bilinear superforms. Moreover, for any decomposition of $`V`$ into the sum of nontrivial graded subspaces $`V=W_1+W_2`$ with nondegenerate restrictions $`f_1`$, $`f_2`$ of $`f`$ one has $`J(V,f)=J(W_1,f_1)+J(W_2,f_2)`$.
Proof. As usual, we identify $`๐ฅ`$ with $`J(V,f)=(F+V_0)+V_1`$ where $`J(V,f)_0=F+V_0`$, $`J(V,f)_1=V_1`$. It follows from the rule of multiplication, that $`๐ฅ_1๐ฅ_1=F1`$, where 1 denotes the identity in $`J(V,f)`$. In particular, $`๐_1๐_1F1`$ and $`_1_1F1`$. Note that the idempotents in $`๐ฅ_0`$ have the form: 1 or $`\frac{1}{2}+v`$ where $`f(v,v)=\frac{1}{4}`$, $`vV_0`$. In particular, if $`v_1`$ and $`v_2`$ are pairwise orthogonal idempotents in $`๐ฅ_0`$, $`v_1=\frac{1}{2}+v`$, $`v_2=\frac{1}{2}v`$ where $`vV_0`$.
Next we consider the following cases.
(a) $`๐K_3`$ where $`K_3=e,x,y`$, $`[x,y]=e`$, $`ex=\frac{x}{2}`$, $`ey=\frac{y}{2}`$, $`e^2=e`$, $`e๐_0`$. As mentioned above, $`e=1`$ or $`e=\frac{1}{2}+v`$, $`vV_0`$. If $`e=1`$, then, obviously, $`ex=x\frac{x}{2}`$. Hence, $`e=\frac{1}{2}+v`$. Then, on the one hand, $`[x,y]=f(x,y)10`$. Conversely, $`[x,y]=\frac{1}{2}+v`$ where $`v0`$. Hence $`\frac{1}{2}+v=f(x,y)1`$. However $`f(x,y)1F`$, $`vF`$, which is wrong.
(b) $`๐D_t`$, $`t1,0,1`$ where $`D_t=e_1,e_2,x,y`$, $`[x,y]=e_1+te_2`$, $`e_1,e_2`$ are pairwise orthogonal idempotents in $`๐_0`$. Hence there exists $`vV_0`$, such that $`e_1=\frac{1}{2}+v`$, $`e_2=\frac{1}{2}v`$. Then $`[x,y]=\frac{1+t}{2}1+(1t)v`$. On the other hand, $`[x,y]=f(x,y)10`$. However $`(1t)v=0`$, that is $`t=1`$, but $`D_1J(V^{},f^{})`$.
(c) $`๐osp(n,m)`$. Since $`๐_02`$, then $`n=1`$, $`m=1`$. It is easy to check that $`osp(1,1)_1osp(1,1)_1`$ cannot be generated by one idempotent.
(d) $`๐M_{n,m}(F)`$. As in the previous case, $`\text{rk}๐_02`$, that is, $`n,m=1`$. By simple calculation we can show that $`M_{1,1}(F)M_{1,1}(F)`$ cannot be a linear span of one idempotent.
(e) $`๐P(n)`$ or $`Q(n)`$. In this case, the proof is similar to one in cases (c) and (d).
It remains to prove the second part of Theorem 5.1, that is, the existence of the decomposition. A $`Z_2`$-graded vector space $`V=V_0+V_1`$ can be represented as the sum of two $`Z_2`$-graded vector subspaces $`W_1`$ and $`W_2`$ in such a way that $`V_0=(W_1)_0+(W_2)_0`$ $`V_1=(W_1)_1+(W_2)_1`$ and the restriction of $`f`$ to $`(W_1)_0`$, $`(W_2)_0`$, $`(W_1)_1`$, $`(W_2)_1`$ are non-singular. Thus, $`๐ฅ=(F+V_0)+V_1=(F+(W_1)_0+(W_2)_0)+((W_1)_1+(W_2)_1)`$ is the sum of two proper simple subsuperalgebras of types $`J(W_1,f_1)`$ and $`J(W_2,f_2)`$, respectively.
Decompositions of $`K_3`$
Let $`๐`$ be a subsuperalgebra of $`K_3`$. Then we have the following restrictions: $`dim๐3`$ and $`\text{rk}๐_0=1`$. Considering all cases one after another, we obtain that $`๐J(V,f)`$ is the only possible case, and $`dim๐=2`$. Hence $`dim๐_0=1`$ $`๐_0=(K_3)_0=e`$, where $`e`$ is an idempotent. In other words, $`e`$ is the identity of $`๐`$. However, if we consider some element of the form $`\alpha x+\beta y`$ belonging to $`๐_1`$, then $`e(\alpha x+\beta y)=\frac{\alpha x+\beta y}{2}`$, that is, $`e`$ cannot be the identity. This implies that $`K_3`$ has no subsuperalgebras of the type $`J(V,f)`$. As a direct consequence, we have
###### Theorem 5.2
A Jordan superalgebra of the type $`K_3`$ has no decompositions into the sum of two proper simple non-trivial subsuperalgebras.
Decompositions of $`D_t`$
Acting in the same manner as in the previous case, we come to the conclusion that the only possible subsuperalgebras in $`D_t`$ can be of types $`K_3`$ and $`J(V,f)`$. Let $`๐K_3`$ be a subsuperalgebra of $`D_t`$. Then we can choose a basis in $`๐`$ in such a way that $`๐=e^{},x^{},y^{}`$, $`e^{}x^{}=\frac{x^{}}{2}`$, $`e^{}y^{}=\frac{y^{}}{2}`$, $`[x^{},y^{}]=e^{}`$, $`e^2=e^{}`$. Moreover, $`๐_0^{}=e^{}`$, $`๐_1^{}=x^{},y^{}`$. Therefore, $`e`$ is an idempotent in $`(D_t)_0=e_1,e_2`$. Hence either $`e^{}=e_i`$, $`i=1,2`$, or $`e^{}=e_1+e_2`$. In the last case, we have $`e^{}(\alpha x+\beta y)=(e_1+e_2)(\alpha x+\beta y)=(\alpha x+\beta y)\frac{(\alpha x+\beta y)}{2}`$. Hence $`e^{}=e_i`$. On the other hand, $`[(\alpha x+\beta y),(\alpha ^{}x+\beta ^{}y)]=(\alpha \beta ^{}\beta \alpha ^{})(e_1+te_2)e^{}=e_i`$. This implies that $`๐`$ of the type $`K_3`$ cannot be a subsuperalgebra $`D_t`$.
Let $`๐J(V,f)`$ be a subsuperalgebra of $`D_t`$. Then $`๐_0(D_t)_0=e_1,e_2`$, $`๐_1(D_t)_1=x,y`$. It is well-known that a subsuperalgebra of type $`J(V,f)`$ has the identity $`e`$, $`ee_1,e_2`$, that is, $`e=e_1+e_2`$. If $`dim๐_0>1`$, then we always can choose some element of form $`(\alpha e_1+\beta e_2)`$ which is linearly independent with $`e_1+e_2`$, and $`(\alpha e_1+\beta e_2)^2`$ is proportionate to $`e_1+e_2`$. This implies that $`\alpha =\beta `$, $`dim๐_0=1`$. However if $`D_t=๐+`$ where $`๐J(V_1,f_1)`$, $`J(V_2,f_2)`$, then $`๐_0=_0=e_1+e_2`$, that is, $`(D_t)_0๐_0+_0`$.
###### Theorem 5.3
A Jordan superalgebra of the type $`D_t`$ has no simple decomposition into the sum of two non-trivial subsuperalgebras.
## 6 Acknowledgment
The author uses this opportunity to thank her supervisor Prof. Bahturin for his helpful cooperation, many useful ideas and suggestions.
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# Differences between application of some basic principles of quantum mechanics on atomic and mesoscopic levels
## I Introduction
Richard Feynman remarked: โI think I can safely say that nobody today understands quantum physicsโ. This remark may seem queer for people who studied and use quantum physics but some experts understand that in contrast to the theories of relativity, quantum mechanics is not yet based on a generally accepted conceptual foundation Nikulov01 . Not only the collision of principles of quantum mechanics with macroscopic realism Nikulov02 ; Nikulov03 and the Einstein-Podolsky-Rosen paradox Nikulov04 ; Nikulov05 are indicative of our incomprehension of quantum physics. There are some quantum effects observed, first of all, on the mesoscopic level, strangeness of which is disregarded by most scientists who do not understand that nobody today understands quantum physics.
The experimental results corroborate for the present all principles of quantum physics Nikulov06 , even in defiance of common sense Nikulov07 ; Nikulov08 . But the essence of these principles is not clear and is discussed now actively Nikulov09 ; Nikulov10 ; Nikulov11 . The collision between quantum mechanics and macroscopic realism Nikulov02 ; Nikulov03 should be expected on the mesoscopic level. Therefore the consideration of differences between application of basic principles of quantum mechanics on atomic and mesoscopic levels is most urgent.
## II Quantum mechanics versus macroscopic realism
One of the three โaxisesโ along which, according to A.J. Leggett Nikulov12 , it is not unreasonable to seek evidence of a breakdown of the quantum mechanics scheme of the physical world is the collision of it with our immediate experience of the โeverydayโ world. The obvious contradiction between the quantum mechanics and macroscopic realism was laid stress by Erwin Schrodinger already seventy years ago Nikulov13 but only in the last years this problem is not only merely philosophical but it can be tested in experiment Nikulov03 first of all on the mesoscopic level, i.e. between the microscopic (atomic) world and the Schrodinger cat. The formalism of the quantum mechanics, its Copenhagen interpretation, was developed first of all for the microscopic (atomic) level and it comes into collision with some logical difficulties on the mesoscopic level.
According to the formalism of the quantum mechanics a quantum system can be in a superposition of states but this superposition can not be observed because of its reduction to single state at measuring. The principle of the impossibility of noninvasive measurement seems admissible on the microscopic level when measuring device can not be smaller than measured object. But we can not assume that the Schrodinger cat can die or revive because of our look. The contradiction between quantum mechanics and the possibility of noninvasive measurability Nikulov02 ; Nikulov14 may can emerge on the mesoscopic level.
## III Quantization of the momentum circulation
Other difficulty can be connected with the quantization of momentum circulation. According to the classical physics the momentum $`p=mv+qA`$ of a particle with a charge $`q`$ should maintain a constant value in absence of any force whereas the quantum number $`n`$ in the relation for the momentum circulation
$$_l๐lp=_l๐l(mv+qA)=m_l๐lv+q\mathrm{\Phi }=n2\pi \mathrm{}$$
$`(1)`$
can change without any evident force. There is not problem on the microscopic realm, where electrons do not change their state of motion in the absence of an electromagnetic force but the problem is on the mesoscopic level Nikulov15 . The mysterious change of state of electron motion without forces acting on the electrons can be both in superconductor Nikulov15 and other (semiconductor and normal metal) mesoscopic structures with the quantization (1) of momentum circulation.
The quantization (1) takes place $`_l๐lp=n2\pi \mathrm{}`$ when the wave function of a particle is closed in a two-connected mesoscopic loop and $`m_l๐lv=n2\pi \mathrm{}q\mathrm{\Phi }=2\pi \mathrm{}(n\mathrm{\Phi }/\mathrm{\Phi }_0)0`$, i.e. the state with zero velocity $`v=0`$ is forbidden, when the magnetic flux $`\mathrm{\Phi }`$ inside the loop is not divisible by the flux quantum $`\mathrm{\Phi }_0=2\pi \mathrm{}/q`$. On the other hand the velocity can be zero $`v=0`$ in the state with unclosed wave function when the quantization (1) is not valid. In this case the circular velocity of the particle $`v`$ should change, i.e. the particle should accelerate, from $`v=0`$ to $`v=_l๐lv/l=2\pi \mathrm{}(n\mathrm{\Phi }/\mathrm{\Phi }_0)/l`$ and the momentum circulation should change from $`q\mathrm{\Phi }`$ to $`n2\pi \mathrm{}`$ at the closing of the wave function without any evident force.
There is important to accentuate a fundamental difference between atomic and mesoscopic levels. A switching between states with different connectivity of wave function can not be realized on atomic level whereas it can be enough easy made on mesoscopic level. For example it can be realized by switching of a segment $`l_s`$ of a loop $`l`$ between superconducting, i.e. with a density of superconducting pairs $`n_s>0`$, and normal states with $`n_s=0`$, whereas other segment $`l_{scs}=ll_s`$ remaining all time in superconducting state with $`n_s>0`$ JLTP98 ; PRB01 . The quantization (1) should be along any closed path $`l`$ of the loop circumference when $`n_s>0`$ along whole loop and the quantization (1) is not valid along $`l`$ when $`n_s=0`$ in the $`l_s`$ segment. The velocity of superconducting pairs
$$_l๐lv_s=\frac{2\pi \mathrm{}}{m}(n\frac{\mathrm{\Phi }}{\mathrm{\Phi }_0})$$
$`(2)`$
and a density of the persistent current $`j_p=2en_sv_s0`$ should be nonzero along $`l`$ in the closed superconducting state at $`\mathrm{\Phi }n\mathrm{\Phi }_0`$ because of the quantization (1), whereas equilibrium velocity $`v_s=0`$ and current $`j_p=0`$ in the $`l_{scs}`$ segment when the $`l_s`$ segment is in the normal state with a non-zero resistance $`R_{ls}>0`$. Thus, superconducting pairs in the $`l_{scs}`$ segment should accelerate without any force, in contradiction with the law of momentum conservation, at the switching of the $`l_s`$ segment from the normal $`n_s=0`$ to superconducting $`n_s>0`$ state. This change can be fixed experimentally by way of an observation of the appearance of the persistent current at closing of superconducting state.
The term โpersistent currentโ was at first used for the current in superconducting state SupPC61 ; SupPC63 ; SupPC65 , i.e. at $`T<T_c`$. Under equilibrium conditions at $`T<T_c`$ the quantization (1) is valid during all time since coherence of wave function of superconducting pairs exists until the superconducting state exists. Above superconducting transition $`T>T_c`$ superconducting pairs exist because of thermal fluctuations Skocpo75 and coherence of their wave function along whole loop $`l`$ appears only at times. It is enough in order the persistent current, i.e. a direct circular current observed under equilibrium conditions, exists not only at $`T<T_c`$ but also in non-superconducting state at $`T>T_c`$ Kulik1 , when the resistance along $`l`$ is not zero $`R_l>0`$. First experimental evidence of the persistent current at $`R_l>0`$ in the fluctuation region $`TT_c`$ is the Little-Parks oscillations of the resistance of cylinder LitPar62 or loop LitPar92 in magnetic field $`R_l(\mathrm{\Phi }/\mathrm{\Phi }_0)`$.
The observation of the circular persistent current $`I_p`$ at a constant magnetic field $`d\mathrm{\Phi }/dt=0`$ in a loop with a non-zero resistance $`R_l>0`$ contradicts to the habitual knowledge according to which such current should disappear without the Faradayโs voltage $`_l๐lE_F=d\mathrm{\Phi }/dt=0`$ because of dissipation, $`R_lI_p0`$, during the time of current relaxation $`\tau _{RL}=L_l/R_l`$. According to the explanation PRB01 the persistent current does not disappear at $`R_l>0`$ since the velocity decrease because of the dissipation force is compensated by the velocity change because of the quantization (1) at closing of superconducting state at reiterate switching of the loop by thermal fluctuations between superconducting states with different connectivity. The explanation PRB01 of the observation of the persistent power $`R_lI_p^20`$ as a fluctuation phenomenon is natural since $`I_p0`$ at $`R_l>0`$ is observed only in the fluctuation region near $`T_c`$, where the loop is switched by fluctuations between superconducting states with different connectivity. According to this explanation PRB01 the observation of the persistent current $`I_p0`$ at $`R_l>0`$ in the fluctuation region of superconducting loop is experimental evidence of violation of the law of conservation of momentum circulation. Already the observation of the direct circular current $`I_p`$ at $`d\mathrm{\Phi }/dt=0`$ and $`R_l>0`$ is challenge to this law since it is observed at $`R_l>0`$, as well as a conventional circular current, but without the circular Faradayโs force $`2eE_F`$, $`_l๐l2eE_F=2ed\mathrm{\Phi }/dt=0`$.
The wave function not only superconducting pairs in the fluctuation region at $`T>T_c`$ but also of electrons in mesoscopic semiconductor and normal metal loops can become closed at times. I.O.Kulik predicted first the persistent current in normal metal mesoscopic structure Kulik2 just after the consideration of this quantum phenomenon at $`T>T_c`$ in superconductor Kulik1 . It is much more difficult to observed the persistent current of electron than superconducting pairs. Nevertheless the advancement of cryogenic and microfabrication technologies had allowed to make attempts to observe the persistent current in semiconductor Semic93 ; Semic95 ; Semic01 and normal metal Normal90 ; Normal91 ; Normal01 nanostructures. First it was made only in 1990, i.e. in 20 years after the predictionKulik2 . It may be therefore most authors refer to Buttiker as the first prediction of the persistent current in non-superconducting structures. The persistent current in non-superconducting loops also contradicts to the habitual knowledge since the resistance of these loop is not zero.
An additional, more obvious, experimental evidence of violation of the law of momentum conservation is the observation Dub02 ; Nikulov16 of the quantum oscillations of the dc voltage $`V_{dc}(\mathrm{\Phi }/\mathrm{\Phi }_0)`$ on segments of asymmetric superconducting loops predicted in JLTP98 ; PRB01 . The potential difference $`R_{ls}I_p`$ should appear on the segment $`l_s`$ just after its switching in the normal state with $`R_{ls}>0`$ if the persistent current in the loop $`l`$ was non-zero $`I_p0`$ before the switching. This potential difference $`V(t)=R_{ls}I(t)=R_{ls}I_p\mathrm{exp}(t/\tau _{RL})`$, as well as the circular current $`I(t)=I_p\mathrm{exp}(t/\tau _{RL})`$, are extinguished during a finite time of current relaxation $`\tau _{RL}=L_l/R_{ls}`$ because of a finite value of the loop inductance $`L_l`$. The time average of the $`V(t)`$ voltage during the time $`t_n`$ of a staying of the $`l_s`$ in the normal state $`\overline{V}^{t_n}=t_n^1_0^{t_n}V(t)=R_{ls}I_pt_n^1_0^{t_n}\mathrm{exp}(t/\tau _{RL})`$ equals $`\overline{V}^{t_n}R_{ls}I_p`$ at $`t_n\tau _{RL}`$ and $`\overline{V}^{t_n}L_lI_p/t_n`$ at $`t_n\tau _{RL}`$. The dc component of the voltage measured during a long time $`\mathrm{\Theta }`$, $`V_{dc}=\mathrm{\Theta }^1_\mathrm{\Theta }๐tV(t)=N_{sw}^1_{Nsw}R_{ls}I_p\omega _{sw}_0^{t_n}\mathrm{exp}(t/\tau _{RL})`$ equals $`V_{dc}\overline{R_{ls}I_pt_n}\omega _{sw}`$ at $`t_n\tau _{RL}`$ and $`V_{dc}L_l\omega _{sw}\overline{I_p}`$ at $`t_n\tau _{RL}`$ in the case of reiterate switching of the $`l_s`$ segment between superconducting and normal states with a frequency $`\omega _{sw}=N_{sw}/\mathrm{\Theta }`$.
The switching of the $`l_s`$ with the frequency $`\omega _{sw}`$ means that during the long time $`\mathrm{\Theta }`$ the loop $`l`$ is $`N_{sw}`$ times in the closed superconducting state and $`N_{sw}`$ times in the unclosed superconducting state. The density of the persistent current $`j_p=2en_sv_s`$ in each (from $`N_{sw}`$) closed superconducting state is determined by the $`n_s`$ value and the quantization of the velocity (2). The density $`j_p`$ is uniform across the narrow section $`s\lambda _L^2`$ of the loops measured in JLTP98 ; PRB01 . Where $`\lambda _L`$ is the London penetration depth. The persistent current in the closed superconducting state of the loop equals $`I_p=sj_p=s2en_sv_s=(2e\pi \mathrm{}/lm<(sn_s)^1>)(n\mathrm{\Phi }/\mathrm{\Phi }_0)`$ because of the quantization (2) and since its value should be uniform along $`l`$ in the stationary state: $`_l๐lv_s=(I_p/2e)_l๐l(sn_s)^1=(I_p/2e)l<(sn_s)^1>`$. The quantum number $`n`$ can be any integer number in the closed superconducting state but with overwhelming probability $`P_n\mathrm{exp}(E_n/k_BT)`$ the loop switches in the permitted state with lowest energy $`E_n`$ since the energy difference $`E_{n+1}E_n`$ between adjacent permitted states is much higher than the thermal energy $`k_BT`$ PRB01 . Therefore the average value $`\overline{n}=N_{sw}^1_{Nsw}n=_nnP_n`$ is close to the integer number corresponding to the lowest $`v_s^2(n\mathrm{\Phi }/\mathrm{\Phi }_0)^2`$ value and $`\overline{I_p}=N_{sw}^1_{Nsw}I_p`$ is not zero at $`\mathrm{\Phi }n\mathrm{\Phi }_0`$ and $`\mathrm{\Phi }(n+0.5)\mathrm{\Phi }_0`$. $`\overline{I_p}=0`$ at $`\mathrm{\Phi }(n+0.5)\mathrm{\Phi }_0`$ since two permitted states, $`n\mathrm{\Phi }/\mathrm{\Phi }_0=1/2`$ and $`n\mathrm{\Phi }/\mathrm{\Phi }_0=1/2`$, with opposite direction of the persistent current $`I_pn\mathrm{\Phi }/\mathrm{\Phi }_0`$ have the same energy $`(n\mathrm{\Phi }/\mathrm{\Phi }_0)^2=(1/2)^2=(1/2)^2`$ and therefore $`\overline{n}\mathrm{\Phi }/\mathrm{\Phi }_0=1/2+(1/2)=0`$ at $`\mathrm{\Phi }(n+0.5)\mathrm{\Phi }_0`$.
Thus, the dc voltage $`V_{dc}\overline{I_p}\overline{n}\mathrm{\Phi }/\mathrm{\Phi }_0`$, sign and value of which are periodical function of the magnetic flux $`V_{dc}(\mathrm{\Phi }/\mathrm{\Phi }_0)`$ should be observed on the $`l_s`$ segment at its reiterate switching between superconducting and normal states. Just such quantum oscillations of the dc voltage $`V_{dc}(\mathrm{\Phi }/\mathrm{\Phi }_0)`$ were observed in Dub02 ; Nikulov16 . There is important that the dc potential difference $`V_{dc}(\mathrm{\Phi }/\mathrm{\Phi }_0)`$ is observed both on the switched segment $`l_s`$ and other one $`l_{scs}=ll_s`$ remaining all time in superconducting state. The latter is possible since the acceleration of pair in the electric field $`\overline{dp/dt}=2e\overline{E_p}=2eV_{dc}/l_{scs}`$ is equilibrated by the momentum change, i.e. by the acceleration in opposite direction JLTP98 ; PRB01 , because of the quantization (1). The momentum circulation $`_l๐lp`$ of superconducting pair with the charge $`q=2e`$ changes from $`2e\mathrm{\Phi }`$ to $`n2\pi \mathrm{}`$ at each closing of the wave function. The average value of this change $`N_{sw}^1_{Nsw}(2\pi \mathrm{}n2e\mathrm{\Phi })=2\pi \mathrm{}(\overline{n}\mathrm{\Phi }/\mathrm{\Phi }_0)`$ depends periodically on magnetic flux as well as the dc voltage $`V_{dc}(\mathrm{\Phi }/\mathrm{\Phi }_0)(\overline{n}\mathrm{\Phi }/\mathrm{\Phi }_0)`$ observed in Dub02 ; Nikulov16 . The observation of the dc voltage on the $`l_{scs}`$ segment remaining all time in superconducting state contradicts to the law of momentum conservation. The quantum oscillation of the dc voltage $`V_{dc}(\mathrm{\Phi }/\mathrm{\Phi }_0)`$ may be expected also in semiconductor and normal metal asymmetric mesoscopic loops.
## IV Intrinsic breach of symmetry
The law of momentum conservation is connected with symmetry of space and the violation of this law at the closing of the wave function can be connected with the intrinsic breach of symmetry. The experimental evidence of the intrinsic breach of symmetry is even more obvious than violation to the law of momentum conservation. It is observed in Dub02 ; Nikulov16 that the potential electric field $`E_p(\mathrm{\Phi }/\mathrm{\Phi }_0)=V(\mathrm{\Phi }/\mathrm{\Phi }_0)`$ has right or left direction which changes periodically with the value $`\mathrm{\Phi }/\mathrm{\Phi }_0`$ of the magnetic flux. For example, if the $`E_p`$ direction is right at $`\mathrm{\Phi }/\mathrm{\Phi }_0=1/4`$ then it is left at $`\mathrm{\Phi }/\mathrm{\Phi }_0=3/4`$ Dub02 ; Nikulov16 . It is very strange that direction of a vector changes with a scalar value. We should ask: โWhy can the dc electric field $`E_p`$ have right direction at $`\mathrm{\Phi }/\mathrm{\Phi }_0=1/4`$ and left one at $`\mathrm{\Phi }/\mathrm{\Phi }_0=3/4`$?โ There can be only answer: โBecause the loop is asymmetric, for example the lower half is more narrow than the upper one, see Fig.4 in Nikulov16 , and the circular persistent current has contra-clockwise direction at $`\mathrm{\Phi }/\mathrm{\Phi }_0=1/4`$ and clockwise one at $`\mathrm{\Phi }/\mathrm{\Phi }_0=3/4`$โ.
It seems self-evident that any direct current has a direction. Nobody doubts that a conventional direct circular current $`I=R_l^1(d\mathrm{\Phi }/dt)`$ (it is in the stationary regime at $`tL_l/R_l`$) induced in a loop with a resistance $`R_l`$ by the Faradayโs voltage $`_l๐lE=d\mathrm{\Phi }/dt`$ has clockwise or contra-clockwise direction and this direction determines right or left direction of the potential electric field $`E_p=V`$ observed on a loop segment $`l_s`$ the resistivity $`R_{ls}/l_s`$ of which differs from the one $`R_l/l`$ along whole loop $`l`$, when $`V=(R_{ls}/l_sR_l/l)l_sI`$. But it is no so obvious for the persistent current existing because of the Bohrโs quantization, as well as stable electron orbit in atom. There is important to accentuate the fundamental difference of the persistent current, as one of the mesoscopic quantum phenomena, from the conventional current, on the one hand, and from electron orbit in atom (1), on the other hand.
The direction of a conventional circular current is determined by the circular Faraday electric field $`_l๐lE=d\mathrm{\Phi }/dt`$. But the persistent current is observed at a constant magnetic flux $`\mathrm{\Phi }`$ and, according to the experimental evidence Dub02 ; Nikulov16 , its direction changes with a scalar value $`\mathrm{\Phi }/\mathrm{\Phi }_0`$ without any external vector factor, i.e. the $`I_p`$ can have different directions at the same direction of the magnetic flux $`\mathrm{\Phi }`$ when the $`\mathrm{\Phi }`$ values are different. The observation Dub02 ; Nikulov16 of a direction of the persistent current is experimental evidence of intrinsic breach of clockwise - counter-clockwise symmetry, since, in contrast to the conventional circular current, the $`I_p`$ direction is not determined by an external vector. The periodical dependence $`I_p(\mathrm{\Phi }/\mathrm{\Phi }_0)V_{dc}(\mathrm{\Phi }/\mathrm{\Phi }_0)`$ of the direction of the persistent current with the period $`\mathrm{\Phi }_0=2\pi \mathrm{}/q`$ is indubitable evidence that this intrinsic breach of symmetry is consequence of the Bohrโs quantization (1).
Bohr postulated the quantization (1), $`_l๐lp=_l๐lmv=n2\pi \mathrm{}`$ at $`\mathrm{\Phi }=0`$, in order to explain the stability of electron orbit in atom. There was a logical difficulty in this model until electron considered as a particle having a velocity $`v`$ since it was impossible to answer on the question: โWhat direction has the velocity of electron on stable atomic orbit?โ The uncertainty relation $`\mathrm{\Delta }p\mathrm{\Delta }l\mathrm{}`$ and the wave quantum mechanics have overcome this difficulty. Electron can not has a certain coordinate on stable atomic orbit with a certain momentum according to the uncertainty relation and therefore it can not have a velocity. It is a wave but not a particle in the case of the Bohrโs quantization on atomic orbit. Therefore the Bohrโs quantization does not break a symmetry on the atomic level. But we see that the breach of symmetry because of the Bohrโs quantization is observed Dub02 ; Nikulov16 on the mesoscopic level.
This intrinsic breach of symmetry is observed since the canonical momentum $`p=mv+qA`$ includes not only velocity $`v`$ but also a magnetic vector potential $`A`$ and therefore sign and value of a circular velocity on the lowest permitted state (2) depend periodidically on the $`\mathrm{\Phi }/\mathrm{\Phi }_0`$. It may be considered as the cause of the periodical changes of equilibrium magnetization $`M(\mathrm{\Phi }/\mathrm{\Phi }_0)`$ of both superconductor and non-superconductor, semiconductor Semic01 and normal metal Normal01 , mesoscopic loops. It is very difficult to investigate experimentally a possibility of like oscillations on atomic level since the Bohrโs radius, a typical atomic size $`r_B0.053nm`$, is much smaller than a radius $`r_B=500nm`$ of the mesoscopic loops. The very high magnetic field $`B>\mathrm{\Phi }_0/\pi r_B^2\mathrm{3\; 10}^5T`$ is needed in order to observe the $`M(\mathrm{\Phi }/\mathrm{\Phi }_0)`$ oscillations on atomic level. It is important to note that the $`M(\mathrm{\Phi }/\mathrm{\Phi }_0)`$ oscillations is challenge to the law of momentum conservation since this periodical change is evidence of change of the quantum number $`n=_l๐lp/2\pi \mathrm{}`$ determining the value of momentum circulation (1). One may assume that this change can be only at a breach of the coherence of wave function along $`l`$.
The intrinsic breach of symmetry on the mesoscopic level because of the Bohrโs quantization is challenge to some basic principle of statistical mechanics and thermodynamics FQMT04 since it violates the postulate of absolute randomness of any equilibrium motion QI2002 .
The work was financially supported by the Presidium of Russian Academy of Sciences in the Program โLow- Dimensional Quantum Structuresโ, by Russian Foundation of Basic Research (Grant 04-02-17068) and by ITCS department of Russian Academy of Sciences in the Program โTechnology Basis of New Computing Methodsโ.
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# Matrix Models
## 1 Introduction
At the beginning let us define the topic of the present lectures. As follows from the title, โMatrix Modelsโ are theories in which the fundamental variable is a matrix. The matrix variable can be a just a constant or a function of time or even be defined as a function over some space-time manifold. With this definition almost any model existing in modern physics e.g. YangโMills theory, theories of Gravity etc., will be a โmatrix theoryโ. Therefore, when speaking on the matrix theory usually a simple structure is assumed, e.g. when fundamental variables are constant or at most time dependent. In the first case, the models of random matrices, one has no time therefore no dynamics. This is a statistical theory describing random matrix distributions. These models are popular in many areas e.g. in the context of description of integrable systems in QCD, or nuclear systems as well as in the study of the lattice Dirac operators (for a review see e.g. Morozov:2005mz ; Verbaarschot:2005rj ; Brody:1981cx ; Guhr:1997ve ; Osborn:1998qb ; Verbaarschot:1994qf ; Shuryak:1992pi and references therein). The special case of interest for us are the YangโMills type matrix models arising in String Theory such like the Ishibashi:1996xs .
Another case of of interest are so called matrix mechanics, i.e. theories of time-evolutive matrices. These models along with the random matrix models are of special interest in String Theory. Thus, the YangโMills type matrix models appear to nonperturbatively describe collective degrees of freedom in string theory called *branes*. Branes are extended objects on which the โnormalโ fundamental strings can end. It was conjectured that including the brane degrees of freedom in the โconventionalโ superstring theories leads to their unification into the *M-theory*, a model giving in its different perturbative regimes all known superstring models. The M-theory is believed to be related to the twelve-dimensional membrane. In the light-cone frame it was conjectured to be described by an YangโMills type matrix mechanics (BFSS matrix model) Banks:1996vh . As we will see in the next section, this as well as the IKKT matrix model can be obtained by quantization/deformation of, respectively, the worldvolume of the membrane and the worldsheet of the string.
As it is by now clear, in this notes we are considering mainly these two models, which sometimes are called โmatrix theoriesโ to underline their fundamental role in string theory.
The plan of these note is as follows. In the next section we give the string motivation and introduce the matrix models as dimensional reductions of supersymmetric YangโMills model. Next, we consider NambuโGoto description of the string and membrane and show that the noncommutative deformation of the respectively, worldsheet or worldvolume leads to IKKT or BFSS matrix models. In the following section we analyze the classical solutions to these matrix models and interpret them as noncommutative gauge models. The fact that these models have a common description in terms of the original matrix model allows one to establish the equivalence relations among them.
## 2 Matrix models of String Theory
### 2.1 Branes and Matrices
A breakthrough in the development of string theory, โthe second string revolutionโ happened when it was observed that in the dynamics of fundamental string on has additional degrees of freedom corresponding to the dynamics at the string ends Polchinski:1995mt (see Aharony:1999ti for a review).
In the open string mode expansion the dynamics at the edge is described by an Abelian gauge field (particle) (for a modern introduction to string theory see e.g. Polchinski:book ; Kiritsis:1997hj ). The corresponding charge of the end of the string is called *ChanโPatton factor*. Allowing a superposition of several, say $`N`$ such factors, which correspond to an โ$`N`$-valentโ string end, gives rise to a nonabelian U($`N`$) super YangโMills gauge field in the effective lagrangian of the open string. This is the so called nine-brane. As it was shown in Polchinski:1995mt string theory allows brane configurations of different other dimensions $`p`$, $`0p+110`$. Depending on the type of the string model they preserve parts of supersymmetry.
So, descending down to the lower dimensional $`p`$-branes one gets the $`p+1`$-dimensional reductions of the ten dimensional Super YangโMills model.
It appears that out of all possibilities only two cases are fundamental, namely, this of $`p=0`$ and $`p=1`$. All other cases can be obtained from either $`p=1`$ or $`p=0`$ by condensation of $`1`$\- or $`0`$-branes into higher dimensional objects.
### 2.2 The IKKT matrix model family
As it follows from the space-time picture, the $`1`$ branes are non-dynamical and, therefore, should be described by a random matrix model which is the reduction of the 10d SYM down to zero dimensions:
$$S_1=\frac{1}{4g^2}tr[X_\mu ,X_\nu ]^2tr\overline{\psi }\gamma ^\mu [X_\mu ,\psi ],$$
(1)
where $`g`$ is some coupling constant depending on SYM coupling $`g_{\mathrm{YM}}`$ and the volume of compactification. Matrices $`X_\mu `$, $`\mu =1,\mathrm{},10`$ are Hermitian $`N\times N`$ matrices, $`\psi `$ is a 10d spinor which has $`N\times N`$ matrix index, $`\gamma ^\mu `$ are 10d Dirac $`\gamma `$-matrices.
From the 10d SYM the matrix model (1) inherits the following symmetries:
* Shifts:
$$X_\mu X_\mu +a_\mu ๐,$$
(2)
where $`a_\mu `$ is a c-number.
* SO(10) rotation symmetry
$$X_\mu \mathrm{\Lambda }_\mu {}_{}{}^{\nu }X_{\nu }^{},$$
(3)
where $`\mathrm{\Lambda }SO(10)`$. This is the consequence of the (euclideanized) Lorenz invariance of the ten dimensional SYM model.
* SU(N) gauge symmetry
$$X_\mu U^1X_\mu U,$$
(4)
where $`USU(N)`$, and this is the remnant of the SYM gauge symmetry invariance.
* Also one has left from the SYM model the supersymmetry invariance:
$`\delta _1X_\mu `$ $`=\overline{ฯต}\gamma _\mu \psi ,`$ (5)
$`\delta _1\psi `$ $`=[X_\mu ,X_\nu ]\gamma ^{\mu \nu }ฯต,`$ (6)
as well as the second one which is simply the shift of the fermion,
$`\delta _2X_\mu `$ $`=0,`$ (7)
$`\delta _2\psi `$ $`=\eta ,`$ (8)
where $`ฯต`$ and $`\eta `$ are the supersymmetry transformation parameters.
###### Exercise 1
Find the relation between the coupling $`g`$ in (1) on one side and SYM coupling $`g_{\mathrm{YM}}`$ and the size/geometry of compactification on the other side.
Hint: Use an appropriate gauge fixing.
###### Exercise 2
Show that (35) are indeed the symmetries of the action (1).
The purely bosonic version of IKKT matrix model can be interpreted as the algebraic version of much older EguchiโKawai model Eguchi:1982nm . The last is formulated in terms of SU(N) *group* valued fields $`U_\mu `$ (in contrast to the algebra valued $`X_\mu `$). The action for the EguchiโKawai model reads as,
$$S_{\mathrm{EK}}=\frac{1}{4g_{\mathrm{EK}}^2}\underset{\mu ,\nu }{}tr(U_\mu U_\nu U_\mu ^1U_\nu ^1๐).$$
(9)
By the substitution, $`U_\mu =\mathrm{exp}aX_\mu `$, $`g_{\mathrm{EK}}^2=g^2a^{4d}`$ and taking the limit $`a0`$ one formally comes to the bosonic part of the IKKT action (1).
Note: From the string interpretation we will discuss in the next section it is worth to add an extra term to the IKKT action (1) and, namely, the chemical potential term,
$$\mathrm{\Delta }S_{chem}=\beta tr๐,$$
(10)
which โcontrolsโ the statistical behavior of $`N`$. In the string/brane picture $`\beta `$ plays the role of the chemical potential for the number of branes. This produces the relative weights for the distributions with different $`N`$, which can not be catched from the arguments we used to write down the action (1).
### 2.3 The BFSS model family
Let us consider another important model which describes the dynamics of zero branes Banks:1996vh . Basic ingredients of this model are roughly the same as for the previous one, the IKKT model, except that now the matrices depend on time. The action for this model is the dimensional reduction of the ten dimensional SYM model down to the only time dimension:
$$S_{\mathrm{BFSS}}=\frac{1}{g_{\mathrm{BFSS}}}dttr\left\{\frac{1}{2}(_0X_i)^2+\overline{\psi }_0\psi \frac{1}{4}[X_i,X_j]^2\overline{\psi }\gamma _i[X_i,\psi ]\right\},$$
(11)
where, now, the index $`i`$ runs from one to nine.
The action (11) describes the dynamics of zero branes in IIA string theory, but it was also proposed as the action for the M-theory membrane in the light-cone approach. As we are going to see in the next section, this model along with the IKKT model can be obtained by worldvolume quantization of the membrane action.
Another known modification of this action for the pp-wave background was proposed by BerensteinโMaldacenaโNastase (BMN) Berenstein:2002jq ; Berenstein:2002zw ; Berenstein:2002sa . It differs from the BFSS model additional terms which are introduced in order to respect the pp-wave supersymmetry. The action of the BMN matrix model reads:
$$\begin{array}{c}S_{\mathrm{BMN}}=\mathrm{d}ttr[\frac{1}{2(2R)}(_0X_i)^2+\overline{\psi }_0\psi \hfill \\ \hfill +\frac{(2R)}{4}[X_i,X_j]^2\mathrm{i}(2R)\overline{\psi }\gamma ^i[X_i,\psi ]]+S_{mass},\end{array}$$
(12)
where $`S_{mass}`$ is given by
$$\begin{array}{c}S_{mass}=\mathrm{d}ttr[\frac{1}{2(2R)}(\left(\frac{\mu }{3}\right)\underset{i=1,2,3}{}X_i^2\left(\frac{\mu }{6}\right)\underset{i=4,\mathrm{},9}{}X_i^2)\hfill \\ \hfill \frac{\mu }{4}\overline{\psi }\gamma _{123}\psi \frac{\mu \mathrm{i}}{3}\underset{jkl=1,\mathrm{},3}{}ฯต_{ijk}X_iX_jX_k]\end{array}$$
(13)
The essential difference of this model from the standard BFSS one is that due to the mass and the Chern-Simons terms this matrix model allow stable vacuum solutions which can be interpreted as spherical branes (see e.g. Valtancoli:2002rx ; Sochichiu:2002ta . Such vacuum configurations can not exist in the original BFSS model.
## 3 Matrix models from the noncommutativity
In this section we show that the Matrix models which we introduced in the previous section arise when one allows the worldsheet of the string/worldvolume of the membrane to possess noncommutativity. It is interesting to note from the beginning that the โquantizationโ of the string worldsheet leads to the IKKT matrix model, while the space noncommutative membrane is described by the BFSS model. Let us remind that the above matrix models were introduced to describe, respectively, the $`1`$\- and 0-branes, while the string and the membrane are respectively 1- and 2-brane objects. In the shed of the next section this can be interpreted as deconstruction of the 1- and 2-branes into their basic components, namely $`1`$\- and 0-brane objects.
In this section we consider only the bosonic parts. The extension to the fermionic part is not difficult, so this is left to the reader as an exercise.
### 3.1 Noncommutative string and the IKKT matrix model
In trying to make the fundamental string noncommutative one immediately meets the following problem: The noncommutativity parameter is a dimensional parameter and, therefore, hardly compatible with the worldsheet conformal symmetry which plays a fundamental role in the string theory. Beyond this there is no theoretical reason to think that the worldsheet of the fundamental string should be noncommutative. On the other hand, the are other string-like objects in the nonperturbative string theory: D1-branes or D-strings. As it was realized, in the presence of the constant nonzero Neveu-Schwarz $`B`$-field the brane can be described by a noncommutative gauge models Cheung:1998nr ; Chu:1998qz ; Chu:1999gi ; Seiberg:1999vs . Then, in contrast to the fundamental string, it is natural to make the D-string noncommutative.
Let us start with the Euclidean NambuโGoto action for the string,
$$S_{\mathrm{NG}}=T\mathrm{d}^2\sigma \sqrt{\underset{ab}{det}_aX^\mu _bX_\mu },$$
(14)
where $`T`$ is the D-string tension and $`X^\mu =X^\mu (\sigma )`$ are the embedding coordinates. The expression under the square root of the r.h.s. of (14) can equivalently rewritten as follows,
$$detXX=\frac{1}{2}\mathrm{\Sigma }^2,$$
(15)
where,
$$\mathrm{\Sigma }^{\mu \nu }=ฯต^{ab}_aX^\mu _bX^\nu ,$$
(16)
which is the induced the worldsheet volume form of the embedding $`X^\mu (\sigma )`$.
The NambuโGoto action then becomes:
$$S_{\mathrm{NG}}=T\mathrm{d}^2\sigma \sqrt{\frac{1}{2}\mathrm{\Sigma }^2}.$$
(17)
This action is nonlinear and still quite complicate. A much simple form can be obtained using the Polyakov trick. To illustrate the idea of te trick which is widely used in the string theory consider first the example of a particle.
#### Polyakovโs trick
The relativistic particle is described by the following reparametrization invariant action,
$$S_p=md\tau \sqrt{\dot{x}^2},$$
(18)
where $`m`$ is the mass and $`x`$ is the particle coordinate. The dynamics of the particle (18) is equivalent, at least classically to one described by the following action,
$$S_{pp}=d\tau \left(\frac{1}{2}e^1\dot{x}^2+m^2e\right).$$
(19)
In this form one has a new variable $`e`$ which plays the role of the line einbein function, or better to say of the one-dimensional volume form.
To see the classical equivalence between (19) and (18) one should write down the equations of motion arising from the variation of $`e`$,
$$e^2=\frac{\dot{x}^2}{m^2},$$
(20)
and use it to substitute $`e`$ in the action (19) which should give exactly (18).
###### Exercise 3
Show this!
As one can see, both actions (18) and (19) are reparametrization invariant, the difference being that the Polyakov action (19) is quadratic in the particle velocity $`\dot{x}`$. This trick is widely used in the analysis of nonlinear systems with gauge symmetry. In what follows we will apply it too.
Let us turn back to our string and the action (17). Applying the Polyakov trick, one can rewrite the action (17) in the following (classically) equivalent form,
$$S_{\mathrm{NGP}}=\mathrm{d}^2\sigma \left(\frac{1}{4}\eta ^1\{X_\mu ,X_\nu \}^2+\eta T^2\right),$$
(21)
where $`\eta `$ is the string โareaโ density and we introduced the Poisson bracket notation,
$$\{X,Y\}=ฯต^{ab}_aX_bY.$$
(22)
It is not very hard to see that the bracket defined by (22) satisfies to all properties a Poisson bracket is supposed to satisfy.
###### Exercise 4
Do it!
Let us note, that the Poisson bracket (22) is not an worldsheet reparametrization invariant quantity. Under the reparametrizations $`\sigma \sigma ^{}(\sigma )`$ it transforms like density rather than scalar the same way as $`\eta `$ is:
$`\{X,Y\}`$ $`det\left({\displaystyle \frac{\sigma ^{}}{\sigma }}\right)\{X,Y\}^{}`$ (23a)
$`\eta `$ $`det\left({\displaystyle \frac{\sigma ^{}}{\sigma }}\right)\eta (\sigma ^{})`$ (23b)
Having two densities one can master a scalar,
$$\{X,Y\}_s=\eta ^1\{X,Y\},$$
(24)
which is invariant. Actually, these two definitions coincide in the gauge $`\eta =1`$, which in some cases may be possible only locally. In terms of the scalar Poisson bracket the action is rewritten in the form as follows
$$S_{\mathrm{NGP}}=\mathrm{d}^2\sigma \eta \left(\frac{1}{4}\{X_\mu ,X_\nu \}^2+T^2\right),$$
(25)
where $`\mathrm{d}^2\sigma \eta `$ is the invariant worldsheet area form.
#### โQuantizationโ
Consider the naive quantization procedure we know from the quantum mechanics. The classical mechanics is described by the canonical classical Poisson bracket,
$$\{p,q\}=1,$$
(26)
and the quantization procedure consists, roughly speaking, in the replacement of the canonical variables $`(p,q)`$ by the operators $`\widehat{p},\widehat{q}`$. At the same time the $`\mathrm{i}\mathrm{}\times `$(Poisson bracket) is replaced by the commutator of the corresponding operators. In particular,
$$\{p,q\}[\widehat{p},\widehat{q}]=\mathrm{i}\mathrm{}.$$
(27)
Afterwards, main task consists in finding the irreducible representation(s) of the obtained algebra<sup>1</sup><sup>1</sup>1In fact, the enveloping algebra rather the Lie algebra itself.. From the undergraduate course of quantum mechanics we know that there are many unitary equivalent ways to do this, e.g. the oscillator basis representation is a good choice.
Under the quantization procedure functions on the phase space are replaced by operators acting on the irreducible representation space of the algebra (27). For these functions and operators one have the correspondence between the tracing and the integration over the phase space with the Liouville measure
$$\frac{\mathrm{d}p\mathrm{d}q}{2\pi \mathrm{}}\mathrm{}tr\mathrm{}$$
(28)
Let us turn to our string model. As in the case of quantum mechanics, under the quantization we mean the replacing the fundamental worldsheet variables $`\sigma ^1`$ and $`\sigma ^2`$ by corresponding operators: $`\widehat{\sigma }^1`$ and $`\widehat{\sigma }^2`$, such that the *invariant* Poisson bracket is replaced by the commutator according to the rule:
$$\{,\}_{\mathrm{PB}}=\mathrm{i}/\theta [,],$$
(29)
where $`\theta `$ is the deformation parameter (noncommutativity). The worldsheet functions are replaced by the operators on the Hilbert space on which $`\widehat{\sigma }^a`$ act irreducibly. As we have two forms of the Poisson bracket the question is wether one should use the density form of the Poisson bracket (22) or the invariant form (24)? The correct choice is the invariant form (24). This is imposed by the fact that the operator commutator is invariant with respect of the choice of basic operator set (in our case it is given by operators $`\widehat{\sigma }^a`$).
Let us note that with the choice of invariant Poisson bracket in (29) the operators $`\widehat{\sigma }^a`$, generally, do not have standard Heisenberg commutation relations. Rather than that, they commute to a nontrivial operator,
$$[\widehat{\sigma }^1,\widehat{\sigma }^2]=\mathrm{i}\theta \widehat{\eta ^1},$$
(30)
where the operator $`\widehat{\eta ^1}`$ corresponds to the inverse density of the string worldsheet area (i.e. its classical limit gives this density). At the same time the trace in the quantum case corresponds to the worldsheet integration with the invariant measure
$$\mathrm{d}^2\sigma \eta [\mathrm{}]2\pi \theta tr[\mathrm{}].$$
(31)
Having the โquantization rulesโ (29) and (31) one is able to write down the noncommutative analog of the NambuโGotoโPolyakov string action (21). It looks as follows,
$$S=\alpha tr\frac{1}{4}[X_\mu ,X_\nu ]^2+\beta tr๐,$$
(32)
where $`\alpha `$ and $`\beta `$ are the couplings of the matrix model. In terms of the string and the deformation parameters they read,
$`\alpha `$ $`={\displaystyle \frac{2\pi }{\theta }},`$ (33)
$`\beta `$ $`={\displaystyle \frac{2\pi T^2}{\theta }}.`$ (34)
After the identification of couplings the model (32) is identic with the IKKT model (1). As a bonus we have obtained the chemical potential (10). As we see from the construction, the dimensionality of matrices depend on the irreducibility representation of the noncommutative algebra. As one can expect from what is familiar in quantum mechanics, the compact worldsheets should lead to *finite-dimensional* representations and thus are described, respectively, by matrices of finite dimensions. There is no exact equivalence between the worldsheet geometry and the matrix description. However, the consistency requires that one should recover the worldsheet geometry in the semi-classical limit ($`\theta 0`$).
Another interesting remark is that in this picture the Heisenberg operator basis correspond to the worldsheet parametrization for which $`\eta `$ is constant. as it is well known such parametrization can exist globally only for the topologically trivial worldsheets. On the other hand, in the algebra of operators acting on a separable infinite dimensional Hilbert space one can always find a Heisenberg operator basis.
#### Example I: Torus
To illustrate the above consider the example of quantization of toric worldsheet. The torus can be described by one complex modulus (or two real moduli). We are not interested here in the possible form of the toric metric, so we can choose the parametrization of the torus for which $`\eta =1`$ and the flat worldsheet coordinates span the range
$$0\sigma ^1<l_1,0\sigma ^2<l_2.$$
(35)
The first problem arises when one tries to quantize variables with the range (35). In spite of the fact that the (invariant) Poisson bracket is canonical the operators $`\widehat{\sigma }^{1,2}`$ can not satisfy the Heisenberg algebra,
$$[\widehat{\sigma }^1,\widehat{\sigma }^2]=\mathrm{i}\theta ,$$
(36)
and have bounded values like in (35) at the same time.
###### Exercise 5
Prove this!
To conciliate the compactness and noncommutativity one should use the compact coordinates $`U_a`$ instead,
$$U_a=\mathrm{exp}2\pi \mathrm{i}\widehat{\sigma }_a/l_a,a=1,2.$$
(37)
The compact coordinates $`U_a`$ satisfy the following (Weyl) commutation relations
$$U_1U_2=qU_2U_1,$$
(38)
where $`q`$ is the toric deformation parameter,
$$q=\mathrm{e}^{2\pi ^2\mathrm{i}\theta /l_1l_2}.$$
(39)
If $`q^N=1`$ for some $`N_+`$, then $`U_a`$ generate an irreducible representation of dimension $`N`$. In this case an arbitrary $`N\times N`$ matrix $``$ can be expanded in powers of $`U_a`$, e.g.
$$=\underset{m,n=0}{\overset{N1}{}}M_{mn}U_1^mU_2^n.$$
(40)
Expansion (40) is in terms of monomials in $`U_1`$ and $`U_2`$ ordered in such a way that all $`U_1`$โs are to the left of all $`U_2`$ one can alternatively use the Weyl functions $`W_{mn}`$ defined as
$$W_{mn}=\mathrm{exp}\left(2\pi \mathrm{i}m\widehat{\sigma }_1/l_1+2\pi \mathrm{i}n\widehat{\sigma }_2/l_2\right),$$
(41)
which differs from the product $`U_1^mU_2^n`$ by a polynomial of lower degree, but is symmetrized in $`\widehat{\sigma }^1`$ and $`\widehat{\sigma }^2`$. Using this expansion in terms of the Weyl functions leads one to the description of matrices in terms of the *Weyl symbols* โ ordinary functions subject to the *star product* algebra. Weyl symbols as well as the star product algebras we are going to consider in the next sections.
As a result we have that quantization of the torus surface leads to the description in terms of $`N\times N`$ matrices where the dimensionality $`N`$ of the matrices depends on the torus moduli.
#### Example II: Fuzzy sphere
Another case of interest is the deformation of the spherical string worldsheet. On the sphere there is no global flat parametrization with $`\eta =1`$. It is convenient to represent the two-sphere worldsheet parameters embedded into the three-dimensional Euclidean space:
$$\sigma _1^2+\sigma _2^2+\sigma _3^2=1,$$
(42)
with the induced metric and volume form $`\eta `$. The (invariant) Poisson bracket is given by the following expression<sup>2</sup><sup>2</sup>2We drop out the subscript of the invariant Poisson bracket since it creates no confusion while it is the only used from now on.:
$$\{\sigma _i,\sigma _j\}=(1/r)ฯต_{ijk}\sigma _k.$$
(43)
Quantization of the Poisson algebra (43) leads to the su(2) Lie algebra commutator,
$$[\widehat{\sigma }_i,\widehat{\sigma }_j]=\mathrm{i}(\theta /r)ฯต_{ijk}\widehat{\sigma }_k,$$
(44)
whose unitary irreducible representations are the well known representations of the su(2) algebra. They are parameterized by the spin of the representation $`J`$. The dimensionality of such representation is $`N=2J+1`$. The two dimensional parameters: the radius of the sphere and the noncommutativity parameter are not independent. They satisfy instead,
$$r^4=\theta ^2J(J+1).$$
(45)
Again, arbitrary $`(2J+1)\times (2J+1)`$ matrix can be expanded in terms of symmetrized monomials in $`\sigma _i`$*noncommutative spherical harmonics*, which are the spherical analogues of the Weyl functions.
Turning back to the action one get exactly the same model as in the previous example with $`N=2J+1`$. As a result we get that independently from which geometry one starts one gets basically the same deformed description. The only meaningful parameter is the dimensionality of the matrix and it depends only on the worldsheet area. This is a manifestation of the universality of the matrix description which we plan to explore in the next sections.
### 3.2 Noncommutative membrane and the BFSS matrix model
Let us consider slightly more complicate example, namely that of the membrane. For the membrane one can write a NambuโGoto action too,
$$S_{NG}=T_m_{\mathrm{\Sigma }_3}\mathrm{d}^3\sigma \sqrt{det_aX^\mu _bX_\mu },$$
(46)
where $`T_m`$ is the membrane tension and $`X`$ are the membrane embedding functions.
In the case when the topology of the worldvolume $`\mathrm{\Sigma }_3`$ is of the type $`\mathrm{\Sigma }_3=I\times _2`$, where $`^1`$ is the time interval $`I=[0,t_0]`$ and $`_2`$ is a two dimensional manifold, one has the freedom to choose the worldsheet parameters $`\sigma ^i`$, $`i=1,2,3`$ in such a way that the time like tangential will be always orthogonal to the space-like tangential,
$$_0X^\mu _aX_\mu =0.$$
(47)
In this case the NambuโGoto action takes the following form
$$S_{NG}=T_md\tau \mathrm{d}^2\sigma \sqrt{\frac{1}{2}\dot{X}^2\mathrm{\Sigma }_{\mu \nu }^2},$$
(48)
where,
$$\mathrm{\Sigma }_{\mu \nu }=ฯต_{ab}_aX_\mu _bX_\nu .$$
(49)
In the complete analogy to the case of the string let us rewrite the NambuโGoto action in the Polyakov form,
$$S_{\mathrm{NGP}}=\mathrm{d}^3\sigma \eta \left[\frac{T_m^2}{2}\dot{X}^2+\frac{1}{4}\{X_\mu ,X_\nu \}^2\right],$$
(50)
where the (invariant) Poisson bracket is defined as
$$\{X,Y\}=\eta ^1ฯต_{ab}_X_bY.$$
(51)
Since we partially fixed the reparametrization gauge invariance by choosing the time direction we have the constraint (47). This leads to the following constraint,
$$\{\dot{X}^\mu ,X_\mu \}=0.$$
(52)
Now, straightforwardly repeating the arguments of the previous subsection one can write down the matrix model action. In the present case the action takes the following form:
$$S_m=dt\left(\beta tr\frac{1}{2}\dot{X}^2+\alpha tr\frac{1}{4}[X_\mu ,X_\nu ]^2\right),$$
(53)
where $`\beta =2\pi T^2/\theta `$ and $`\alpha =2\pi /\theta `$, respectively. The action (53) should be supplemented with the following constraint:
$$[\dot{X}_\mu ,X_\mu ]=0.$$
(54)
The constraint (54) can be added to the action (53) with the Lagrange multiplier $`A_0`$. In this case the action acquires the following form:
$$S_{gi}=dt\left(\beta tr\frac{1}{2}(_0X_\mu )^2+\alpha tr\frac{1}{4}[X_\mu ,X_\nu ]^2\right),$$
(55)
which is identic (upto definition of parameters $`\alpha `$ and $`\beta `$) to the bosonic part of the BFSS action (11). By the redefinition of the matrix fields and rescaling of the time one can eliminate the constants $`\alpha `$ and $`\beta `$, so in what follows we can put both to unity.
So far we have considered only the bosonic parts of the membrane. Including the fermions (when they exist) introduces no conceptual changes. Therefore, derivation of the fermionic parts of the IKKT and BFSS matrix model description of the string and membrane is entirely left to the reader.
###### Exercise 6
Derive the fermionic part of both matrix models starting from the superstring/supermembrane.
## 4 Equations of motion. Classical solutions
In this section we consider two types of theories, namely the string and the membrane in the NambuโGoto-Polyakov form and the corresponding matrix models. One can write down equations of motion and try to find out some simple classical solutions in order to compare these cases among each other.
The static equations of motion in the membrane case coincide with the string equations of motion. Therefore, it is enough to consider only the last case: Any solution in the IKKT model has also the interpretation as a classical vacuum of the BFSS theory.
### 4.1 Equations of motion before deformation: NambuโGotoโPolyakov string
Consider first the equations of motion corresponding to the NambuโGotoโPolyakov string (21) in the form one gets just before the deformation procedure.
Variation of $`X_\nu `$ produces the following equations,
$$\{X_\mu ,\eta ^1\{X_\mu ,X_\nu \}\}=0,$$
(56a)
while the variation of $`\eta `$ produces the constraint
$$\eta ^2=\frac{1}{4}\frac{\{X_\mu ,X_\nu \}^2}{T^2}.$$
(56b)
(As in the Polyakov particle case the last equation can be used to eliminate $`\eta `$ from the action (21) in order to get the original NambuโGoto action (14).)
The equations of motion (56) posses a large symmetry related to the reparametrization invariance (23). In order to find some solutions it is useful (but not necessary!) to fix this gauge invariance. As the use of the model is to describe branes, one may be interested in solutions corresponding to infinitely extended branes, which have the topology of $`^2`$. In this, simplest case one can impose the gauge $`\eta =1/4T^2`$. Then, the equations of motion (56) are reduced to
$$\{X_\mu ,\{X_\mu ,X_\nu \}\}=0,\{X_\mu ,X_\nu \}^2=1.$$
(57)
In the case of two dimensions ($`\mu ,\nu =1,2`$), one can find even the generic solution. It is given by an arbitrary canonical transformation $`X_{1,2}=X_{1,2}(\sigma _1,\sigma _2)`$. This is easy to see if to observe that the second equation in (57) requires that the $`XX`$ Poisson bracket must be a canonical one. The first equation is then satisfied automatically. One can also see that all the arbitrariness in the solution is due to the remnant of the reparametrization invariance which is given by the *area preserving diffeomorphisms*. This situation is similar to one met in the case of two dimensional gauge theories where there are no physical degrees of freedom left to the gauge fields beyond the gauge arbitrariness. As we will see later, this similarity is not accidental, in some sense the above matrix model is indeed a two-dimensional gauge theory.
The situation is different in more than two dimensions. In this case we are not able to write down the generic solution, but one can find a significant particular one. The simplest solutions of (57) can be obtained by just lifting up the two-dimensional ones to higher dimensions. In particular, one has the following solution
$$X_1=\sigma _1,X_2=\sigma _2,X_i=0,i=3,\mathrm{},10.$$
(58)
It is not difficult to check that the solution (58) satisfy to both equations (57). The physical meaning of this solution is an infinite Euclidean brane extended in the plane (1,2).
One can see, that by the nature of the model in which fields $`X_\mu `$ are functions of a two dimensional parameter the solutions to the equations of motion are forced always to describe two dimensional surfaces i.e. single brane configurations. One can go slightly beyond this limitation allowing $`X`$โs to be multivalent functions of $`\sigma `$โs. In this case one is able to describe a certain set of multibrane systems, each sheet of $`X`$ corresponding to an individual brane. This situation in application to spherical branes was analyzed in more details in Sochichiu:2002ta .
Another question one may ask is whether one can find solutions describing a compact worldsheet. We are not going to give any proof of the fact that such type of solutions do not exist. Rather we consider a simple example of a cylindrical configuration and show that thew equations of motion are not satisfied by it. An infinite cylinder as an extremal case of the torus can be given by the following parametric description:
$$X_1=\mathrm{sin}\sigma _1,X_2=\mathrm{cos}\sigma _1,X_3=\sigma _2.$$
(59)
The eq. (59) describes a cylinder obtained from moving the circle in the plane (1,2) along the axe 3. The parametrization (59) satisfy the constraint (47), therefore to see wether such surface is a classically stable it is enough to check the the first equation of (57). The explicit evaluation of the equations of motion gives
$`\{X_\mu ,\{X_\mu ,X_1\}\}`$ $`=X_10,`$ (60)
$`\{X_\mu ,\{X_\mu ,X_2\}\}`$ $`=X_20,`$ (61)
$`\{X_\mu ,\{X_\mu ,X_0\}\}`$ $`=0.`$ (62)
As one see, only the equation of motion for the third noncompact direction is satisfied. Other equations can be satisfied if one modifies the action of the model by adding mass terms for e.g. $`X_1`$ and $`X_2`$:
$$SS+m^2(X_1^2+X_2^2).$$
(63)
###### Exercise 7
Modify the classical action in a way to allow the spherical brane solutions. Worldsheet quantize this model and compare it to the BMN matrix model.
Another interesting type of solutions is given by singular configurations with trivial Poisson bracket,
$$\{X_\mu ,X_\nu \}=0.$$
(64)
Obviously, these configurations satisfy the equations of motion. This solution corresponds to an arbitrary open or closed smooth one-dimensional line embedded in $`^D`$. The problem appears when one tries to make this type of solution to satisfy the constraint (47) arising from the gauge fixing $`\eta ^2=1/4T^2`$. This configuration, however is still an acceptable solution before the gauge fixing. The degeneracy of the two dimensional surface into the line results into the degeneracy of the two-dimensional surface reparametrization symmetry into the subgroup of the line reparametrizations. This means in particular that $`\eta ^2=1/4T^2`$ is not an acceptable gauge condition in this point, one must impose $`\eta =0`$ instead.
Let us now turn to the noncommutative case and see how the situation is changed there.
### 4.2 Equations of motion after deformation: <br>IKKT/BFSS matrix models
After quantization of the worldsheet/worldvolume we are left with no Polyakov auxiliary field $`\eta `$. The role of this field in the noncommutative theory is played by the choice of the representation. As most cases we can not smoothly variate the representation, we have no equations of motion corresponding to this parameter. So, we are left with only equations of motion corresponding to the variation of $`X`$โs. For the IKKT model these equations read
$$[X_\mu ,[X_\mu ,X_\nu ]]=0,$$
(65)
while for the BFSS model the variation of $`X`$ leads to the following dynamical equations,
$$\ddot{X}_\mu +[X_\mu ,[X_\mu ,X_\nu ]]=0,$$
(66)
where we also put the brane tension to unity: $`T=1`$. If one is interested in only the static solutions ($`\dot{X}=0`$) to the BFSS equations of motion, then the equation (66) is reduced down to the IKKT equation of motion. Therefore, in what follows we consider only the last one.
By the first look at the equation (65) it is clear that one can generalize the string soluiton (58) from the commutative case. Namely, one can check that the configuration
$$X_1=\widehat{\sigma }_1,X_2=\widehat{\sigma }_2,X_i=0,i=3,\mathrm{},D,$$
(67)
satisfy the equations of motion (65). By the analogy with the commutative case we can say that this configuration describes either Euclidean D-string (IKKT) or a static membrane (BFSS). The solution (67) corresponds to the Heisenberg algebra
$$[X_1,X_2]=1,$$
(68)
which allows only the infinite-dimensional representation. The value of $`X`$ are not bounded, therefore this solution corresponds to a noncompact brane.
What is the role of the $`\eta `$-constraint here? The algebra (68) does not completely specify the solution unless the nature of its representation is also given. In particular, the algebra of $`\widehat{\sigma }`$โs can be irreducibly represented on the whole Hilbert space. In the semiclassical limit this can be seen to correspond to the constraint of the previous subsection.
As we discussed in the case of commutative string, any solution to the equations of motion describes a two dimensional surface and, therefore, has the Poisson bracket of the rank (in indices $`\mu `$ and $`\nu `$) two or zero. In contrast to this, in the noncommutative case one may have solutions with an arbitrary even rank between zero and $`D`$. Indeed, consider a configuration,
$$X_a=p_a,a=1,\mathrm{},p+1,X_i=0,i=p+2,\mathrm{},D,$$
(69)
such that
$$[p_a,p_b]=\mathrm{i}B_{ab},detB0,$$
(70)
where $`B`$ is the matrix with c-number entries $`B_{ab}`$. Such set of operators always exists if the Hilbert space is infinite dimensional separable. The set of operators $`p_a`$ generate a Heisenberg algebra. Interesting cases are when the Heisenberg algebra (70) is represented irreducibly on the Hilbert space of the model, or when this irreducible representation is $`n`$-tuple degenerate. Analysis of these cases we will do in the next sections.
How about the compact branes? As we have already discussed in the previous section, the compact worldsheet solution corresponds to finite dimensional matrices $`X_\mu `$. As it appears for such matrices the only solution to the equation of motion which exists is one with the trivial commutator,
$$[X_\mu ,X_\nu ]=0.$$
(71)
To prove this fact, suppose we find such a solution with $`B_{\mu \nu }=[X_\mu ^{(0)},X_\nu ^{(0)}]0`$ and satisfying the equations of motion (65). The IKKT action (BFSS energy) computed on such a solution is
$$S(X)=\frac{1}{4}B_{\mu \nu }^2tr๐0.$$
(72)
Since this is a solution to the equations of motion the variation of the action should vanish on the solution,
$$\delta S=tr\frac{\delta S}{\delta X_\mu }(X^{(0)})\delta X_\mu =0,\text{for }\delta X_i,$$
(73)
which is not the case: Take $`\delta X_\mu =ฯตX_\mu ^{(0)}`$ to find out that $`\delta S|_{X^{(0)}}0`$. So there are no solutions with nontrivial commutator for the finite dimensional matrix space.
Consider now the extremal case of singular solutions with vanishing commutators,
$$[X_\mu ,X_\nu ]=0.$$
(74)
Obviously, from the equation (74) automatically follows that the equations are satisfied too. This solution exists in both finite as well as infinite-dimensional cases. Since the commutativity of $`X_\mu `$โs allows their simultaneous diagonalization
$$X_\mu =\left(\begin{array}{cccc}x_1^\mu & & & \\ & x_2^\mu & & \\ & & \mathrm{}& \end{array}\right),$$
(75)
this means that the branes which are described by the matrix models are localized $`x_k^\mu `$ being the coordinates of the $`k`$-th brane.
#### The symmetry of the solutions
The various types of solutions have different symmetry properties. Thus, the solution of the type (69) with the algebra of $`p_a`$โs irreducibly represented over the Hilbert space of the model has no internal symmetries. Indeed, by the Schurrโs lemma any operator commuting with all $`p_a`$ is proportional to the identity. In the case when the representation is $`n`$-tuple degenerate one has an U$`(n)`$ symmetry mixing the representations. The degenerate case (74), when $`B_{\mu \nu }=0`$ give rise to some symmetries too. Indeed, an arbitrary diagonal matrix commute with all $`X_\mu `$ given by (75). If no two branes are in the same place: $`x_m^\mu x_n^\mu `$ for any $`mn`$, then the configuration breaks the U$`(N)`$ symmetry group (in the finite-dimensional case) down to the the Abelian subgroup U(1)<sup>N</sup>.
## 5 From the Matrix Theory to Noncommutative YangโMills
This and the following section is mainly based on the papers Sochichiu:2000ud ; Sochichiu:2000fs ; Sochichiu:2000bg ; Sochichiu:2000kr ; Sochichiu:2000kz ; Kiritsis:2002py , the reader is also referred to the lecture notes Sochichiu:2002jh and references therein.
The main idea is to use the solutions from the previous section both as classical vacua, such that arbitrary matrix configuration is regarded as a perturbation of this vacuum configuration, and as a basic set of operators in terms of which the above perturbations are expanded. Now follow the details.
### 5.1 Zero commutator case: gauge group of diffeomorphisms
Consider first the case of the solution with the vanishing commutator (74). We are interested in configurations in which the branes form a $`p`$-dimensional lattice. Using the rotational symmetry of the model, one can choose this lattice to be extended in the dimensions $`1,\mathrm{},p`$:
$$X_ap_a,a+1,\mathrm{},p;X_I=0,I=p+1,\mathrm{}D.$$
(76)
Then an *arbitrary* configuration can be represented as
$$X_a=p_a+A_a,X_I=\mathrm{\Phi }_I.$$
(77)
Let us take the limit $`N\mathrm{}`$ and take such a distribution of the branes in which they form an infinite regular $`p`$-dimensional lattice:
$$p_a\lambda n_a,n_a,$$
(78)
such that the Hilbert space can be split in the product of $`p`$ infinite-dimensional subspaces $`_a`$
$$=_{a=1}^p_a,$$
(79)
such that each eigenvalue $`\lambda n_a`$ is non-degenerate in $`_a`$. In this case the operators $`p_a`$ can be regarded as ($`\mathrm{i}`$ times) partial derivatives on a $`p`$-dimensional torus of the size $`1/\lambda `$,
$$p_a=\mathrm{i}_a.$$
(80)
Now let us turn to the perturbation of the vacuum configuration (77) and try to write it in terms of operators $`p_a`$. Since the algebra of $`p_a`$โs is commutative, they alone fail to generate an irreducible representation in terms of which one can expand an arbitrary operator acting on the Hilbert space $``$. One must instead supplement this set with with $`p`$ other operators $`x^a`$, which together with $`p_a`$ form a Heisenberg algebra irreducibly represented on $``$,
$$[x^a,x^b]=0,[p_a,x^b]=\mathrm{i}\delta _a{}_{}{}^{b}.$$
(81)
From the algebra (81) follows that the operators $`x^a`$ have a continuous spectrum which is bounded: $`\pi /\lambda x^a<\pi /\lambda `$. This precisely means that $`x^a`$ are operators of coordinates on the $`p`$-dimensional torus. Then, an arbitrary matrix $`X`$ can be represented as a an operator function of the operators $`p_a`$ and $`x^a`$,
$$X=\widehat{X}(\widehat{p},\widehat{x}).$$
In the โ$`x`$-pictureโ this will be a differential operator $`X(\mathrm{i},x)`$. There are many ways to represent a particular operator $`X`$ as a operator function of $`p_a`$ and $`x^a`$ which is related to the *ordering*. The *Weyl ordering* we will consider in the next subsection, here let us use a different one in which all operators $`p_a`$ are on the right to all $`x^a`$. In such an ordering prescription one can write down a Fourier expansion of the operator in the following form
$$X=\frac{1}{(2\pi )^p}\mathrm{d}^pz\stackrel{~}{X}(z,x)\mathrm{e}^{\mathrm{i}\widehat{p}z}.$$
(82)
In this parametrization the product of two operators is given by an involution product of the symbols:
$$\stackrel{~}{XY}(z,x)=\stackrel{~}{X}\stackrel{~}{Y}(z,x)=\frac{1}{(2\pi )^p}\mathrm{d}^py\stackrel{~}{X}(y,x)\stackrel{~}{Y}(zy,x+y).$$
(83)
The trace of an operator can be computed in a standard way, namely
$$trX=\mathrm{d}^pxx\left|X\right|x=\mathrm{d}^px\stackrel{~}{X}(0,x)=\mathrm{d}^px\mathrm{d}^plX(l,x),$$
(84)
where in the last part $`X(l,x)`$ is the normal symbol of which is obtained by the replacement of operator $`\widehat{p_a}`$ by an ordinary variable $`l_a`$ in the definition (82),
$`X(l,x)`$ $`={\displaystyle \frac{1}{(2\pi )^p}}{\displaystyle \mathrm{d}^pz\stackrel{~}{X}(z,x)\mathrm{e}^{\mathrm{i}lz}},`$ (85)
$`\stackrel{~}{X}(z,x)`$ $`=tr\mathrm{e}^{\mathrm{i}\widehat{p}z}X.`$ (86)
Now we are ready to write down the whole matrix action (32) in terms of the normal symbols. It looks as follows,
$$S=\mathrm{d}^pl\mathrm{d}^px\left(\frac{1}{4}_{ab}^2+\frac{1}{2}(_a\mathrm{\Phi }_I)^2\frac{1}{4}[\mathrm{\Phi }_I,\mathrm{\Phi }_J]_{}^2\right),$$
(87)
where
$`_{ab}(l,x)`$ $`=_aA_b(l,x)_bA(l,x)[A_a,A_b]_{}(l,x),`$ (88)
$`_a\mathrm{\Phi }`$ $`=_a\mathrm{\Phi }+[A_a,\mathrm{\Phi }]_{}(l,x),`$ (89)
$`[A,B]_{}(l,x)`$ $`=AB(l,x)BA(l,x)`$ (90)
and the star product is defined as in (83).
The model defined by the action (87) has the meaning of YangโMills theory with the infinite dimensional gauge group of diffeomorphism transformations generated by the operators
$$T_f=\mathrm{i}f^a(x)_a.$$
(91)
Because of the noncommutative nature of the products involved in the action (87) the local gauge group is not commutative. However, if one tries to write down the group of global gauge symmetry, one finds out that this group is, in fact nothing else that U(1). Changing only slightly the character of the solution one can also get a non-Abelian global group. Indeed, consider the solution as in (76) with the exception that the Hilbert space is not just (79), but is given by the product of parts $`_a`$ at some (positive integer) power $`n`$:
$$=\left(_{a=1}^p_a\right)^n.$$
(92)
Repeating with this solution the same manipulations which lead us to (87) with the only exception that in this case an arbitrary matrix is represented by a $`(n\times n)`$-matrix valued function instead of just โordinaryโ one, we arrive to the action similar to (87) with the exception that the fields take their value in the u$`(n)`$ algebra and the global gauge group is, respectively, U$`(n)`$. We hope, that the things will clarify a lot when the reader will pass the next subsection.
#### Ordinary gauge model?
A question one may ask oneself is if the fluctuations of the matrix models can be restricted in such a way to get a โnormalโ YangโMills theory with a compact Lie group. In the present case one may restrict the fluctuations around the background (76) to depend on $`\widehat{x}^a`$ operators only. This aim can be achieved by imposing the following constraints on the matrices $`X_\mu `$:
$$[x_a,X_b]=\mathrm{i}\delta _{ab},[x_a,X_I]=0.$$
(93)
Let us note that $`X_a`$ and $`x_a`$ do not form the Heisenberg algebra because the commutator between $`X_a`$ do not necessarily vanish:
$$[X_a,X_b]F_{ab}0.$$
(94)
Dynamically, the constraint (93) can be implemented through the modification of the matrix action by the addition of the constraint (93) with the Lagrange multiplier. The modified matrix model action reads:
$$S_c=tr\left(\frac{1}{4}[X_\mu ,X_\nu ]^2+\rho _{\mu \nu }([x_\mu ,X_\nu ]\mathrm{\Delta }_{\mu \nu })+T^2\right),$$
(95)
where $`\rho _{\mu \nu }`$ are the Lagrange multipliers, $`x_\mu =(x_a,0)`$ and $`\mathrm{\Delta }_{\mu \nu }`$ is equal to $`\delta _{ab}`$ when $`(\mu \nu )=(ab)`$ and zero otherwise. The limit $`N\mathrm{}`$ of the matrix model specified by the action (95) produces the Abelian gauge model. Under similar setup one can obtain also nonabelian gauge models.
### 5.2 Nonzero commutator: Noncommutative YangโMills model
In this subsection we consider the matrix action as a perturbation of the background configuration given by (69) and (70). Here we plan to give a more detailed approach also partly justifying the result of the previous subsection. The operators $`p_a`$ generate a $`(p+1)/2`$-dimensional Heisenberg algebra. If this algebra is represented irreducibly on the Hilbert space of the model (which is in fact our choice), then an arbitrary operator acting on this space can be represented as an operator function of $`p_a`$. Let us consider this situation in more details.
Irreducibility of the representation in particular means that any operator commuting with all $`p_a`$ is a $`c`$-number constant. From this follows that the operators
$$P_a=[p_a,],$$
(96)
which are Hermitian on the space of square trace operators equipped with the scalar product $`(A,B)=trA^{}B`$, are diagonalizable and have non-degenerate eigenvalues.
###### Exercise 8
Prove this!
By a direct check one can verify that the operator $`\mathrm{e}^{\mathrm{i}k_a\widehat{x}^a}`$, where $`\widehat{x}^a=\theta ^{ab}\widehat{p}_b`$, $`\theta B^1`$ is an eigenvector for $`P_a`$ with the eigenvalue $`k_a`$:
$$P_a\mathrm{e}^{\mathrm{i}k\widehat{x}}=[p_a,\mathrm{e}^{\mathrm{i}k\widehat{x}}]=k_a\mathrm{e}^{\mathrm{i}k\widehat{x}}.$$
(97)
This set of eigenvectors form an orthogonal basis ($`P_a`$โs are Hermitian). One can normalize the eigenvectors to delta function trace,
$$E_k=c_k\mathrm{e}^{\mathrm{i}k\widehat{x}},trE_k^{}^{}E_k=\delta (k^{}k).$$
(98)
The normalizing coefficients $`c_k`$ can be found from evaluating explicitly the trace of $`\mathrm{e}^{\mathrm{i}(kk^{})\widehat{x}}`$ in (98) and equating it to the Dirac delta. Let us compute this trace and find the respective quotients. To do this, consider the basis where the set of operators $`x^\mu `$ splits in pairs $`p_i`$, $`q^i`$ satisfying the standard commutation relations: $`[p_i,q_j]=\mathrm{i}\theta \delta _{ij}`$.
As we know from courses of Quantum Mechanics the trace of the operator
$$\mathrm{e}^{\mathrm{i}k^{}\widehat{x}}\mathrm{e}^{\mathrm{i}k\widehat{x}}=\mathrm{e}^{\mathrm{i}(kk^{})\widehat{x}}\mathrm{e}^{\frac{\mathrm{i}}{2}k^{}\times k},$$
(99)
can be computed in $`q`$-representation as,
$$tr\mathrm{e}^{\mathrm{i}(kk^{})\widehat{x}}\mathrm{e}^{\frac{\mathrm{i}}{2}k^{}\times k}=dqq\left|\mathrm{e}^{\mathrm{i}(l_i^{}l_i)q^i+(z^iz^i)p_i}\right|q=1/|c_k|^2\delta (k^{}k),$$
(100)
where $`|q`$ is the basis of eigenvectors of $`๐ช^i`$,
$$๐ช^i|q=q^i|q,q^{}|q=\delta (q^{}q),$$
(101)
and $`l_i`$, $`z^i`$ ($`l_i`$, $`z^i`$) are components of $`k_\mu `$ ($`k_\mu ^{}`$) in the in the parameterizations: $`x^\mu p_i,q^i`$. Explicit computation gives,
$$1/|c_k|^2=\frac{(2\pi )^{\frac{p}{2}}}{\sqrt{det\theta }}.$$
(102)
Now, we have the basis of eigenvectors $`E_k`$ and can write any operator $`F`$ in terms of this basis,
$$\widehat{F}=dk\stackrel{~}{F}(k)\mathrm{e}^{\mathrm{i}k\widehat{x}},$$
(103)
where the โcoordinateโ $`\stackrel{~}{F}(k)`$ is given by,
$$\stackrel{~}{F}(k)=\frac{\sqrt{det\theta }}{(2\pi )^{\frac{p}{2}}}tr(\mathrm{e}^{\mathrm{i}k\widehat{x}}\widehat{F}).$$
(104)
Function $`\stackrel{~}{F}(k)`$ can be interpreted as the Fourier transform of a $`L^2`$ function $`F(x)`$,
$$F(x)=dk\stackrel{~}{F}(k)\mathrm{e}^{\mathrm{i}k_\mu x^\mu }=\sqrt{det\theta }\frac{\mathrm{d}k}{(2\pi )^{p/2}}\mathrm{e}^{\mathrm{i}kx}tr\mathrm{e}^{\mathrm{i}k\widehat{x}}\widehat{F}.$$
(105)
And viceversa, to any $`L^2`$ function $`F(x)`$ from one can put into correspondence an $`L^2`$ operator $`\widehat{F}`$ by inverse formula,
$$\widehat{F}=\frac{\mathrm{d}x}{(2\pi )^{p/2}}\frac{\mathrm{d}k}{(2\pi )^{p/2}}F(x)\mathrm{e}^{\mathrm{i}k(\widehat{x}x)}.$$
(106)
Equations (105) and (106) providing a one-to-one correspondence between $`L^2`$ functions and operators with finite trace,
$$tr๐
^{}๐
<\mathrm{},$$
(107)
give in fact formula for the Weyl symbols. By introducing distributions over this space of operators one can extend the above map to operators with unbounded trace.
###### Exercise 9
Check that (105) and (106) lead in terms of distributions to the correct Weyl ordering prescription for polynomial functions of $`p_\mu `$.
Let us note, that the map (105) and (106) can be rewritten in the following form,
$$F(x)=(2\pi )^{p/2}\sqrt{det\theta }tr\widehat{\delta }(\widehat{x}x)\widehat{F},\widehat{F}=\mathrm{d}^px\widehat{\delta }(\widehat{x}x)F(x),$$
(108)
where we introduced the operator,
$$\widehat{\delta }(\widehat{x}x)=\frac{\mathrm{d}^pk}{(2\pi )^p}\mathrm{e}^{\mathrm{i}k(\widehat{x}x)}.$$
(109)
This operator satisfy the following properties,
$`{\displaystyle \mathrm{d}^px\widehat{\delta }(\widehat{x}x)}=๐,`$ (110a)
$`(2\pi )^{p/2}\sqrt{det\theta }tr\widehat{\delta }(\widehat{x}x)=1,`$ (110b)
$`(2\pi )^{p/2}\sqrt{det\theta }tr\widehat{\delta }(\widehat{x}x)\widehat{\delta }(\widehat{x}y)=\delta (xy),`$ (110c)
where in the r.h.s. of last equation is the ordinary delta function. Also, operators $`\widehat{\delta }(\widehat{x}x)`$ for all $`x`$ form a complete set of operators,
$$[\widehat{\delta }(\widehat{x}x),๐
]0F\mathrm{i}.$$
(110d)
The commutation relations of $`\widehat{x}^\mu `$ also imply that $`\widehat{\delta }(\widehat{x}x)`$ should satisfy,
$$[\widehat{x}^\mu ,\widehat{\delta }(\widehat{x}x)]=\mathrm{i}\theta ^{\mu \nu }_\nu \widehat{\delta }(\widehat{x}x).$$
(110e)
In fact one can define alternatively the noncommutative plane starting from operator $`\widehat{\delta }(\widehat{x}x)`$ satisfying (110), with $`\widehat{x}^\mu `$ defined by,
$$\widehat{x}^\mu =\mathrm{d}^pxx^\mu \widehat{\delta }(\widehat{x}x).$$
(111)
In this case (110e) provides that $`\widehat{x}^\mu `$ satisfy the Heisenberg algebra (69), while the property (110d) provides that they form a complete set of operators. Relaxing these properties allows one to introduce a more general noncommutative spaces.
Let us the operator $`\widehat{\delta }(x)`$ in the simplest case of two-dimensional noncommutative plane. The most convenient is to find its matrix elements $`D_{mn}(x)`$ in the oscillator basis given by,
$$|n=\frac{(\widehat{a}^{})^n}{\sqrt{n!}}|0,\widehat{a}|0=0,$$
(112)
where the oscillator operators $`\widehat{a}`$ and $`\widehat{a}^{}`$ are the noncommutative analogues of the complex coordinates,
$$\widehat{a}=\sqrt{\frac{1}{2\theta }}(\widehat{x}^1+\mathrm{i}\widehat{x}^2),\widehat{a}^{}=\sqrt{\frac{1}{2\theta }}(\widehat{x}^1\mathrm{i}\widehat{x}^2);[\widehat{a},\widehat{a}^{}]=1.$$
(113)
Then the matrix elements read
$$D_{mn}(x)=m\left|\widehat{\delta }^{(2)}(\widehat{a}z)\right|n=tr\widehat{\delta }^{(2)}(\widehat{a}z)P_{nm},$$
(114)
where $`P_{nm}=|nm|`$.
As one can see, up to a Hermitian transposition the matrix elements of $`\widehat{\delta }(\widehat{x}x)`$ correspond to the Weyl symbols of operators like $`|mn|`$, or so called Wigner functions. The computation of (114) gives,<sup>3</sup><sup>3</sup>3For the details of computation see e.g. Harvey:2001yn .
$$D_{mn}^\theta (z,\overline{z})=(1)^n\left(\frac{2}{\sqrt{\theta }}\right)^{mn+1}\sqrt{\frac{n!}{m!}}\mathrm{e}^{z\overline{z}/\theta }\left(\frac{z^m}{\overline{z}^n}\right)L_n^{mn}(2z\overline{z}/\theta ),$$
(115)
where $`L_n^{mn}(x)`$ are Laguerre polynomials,
$$L_n^\alpha (x)=\frac{x^\alpha \mathrm{e}^x}{n!}\left(\frac{\mathrm{d}}{\mathrm{d}x}\right)^n(\mathrm{e}^xx^{\alpha +n}).$$
(116)
It is worthwhile to note that in spite of its singular origin the symbol of the delta operator is a smooth function which is rapidly vanishing at infinity. The smoothness comes from the fact that the operator elements are written in an $`L^2`$ basis. In a non-$`L^2`$ basis, e.g. in the basis of $`x_1`$ eigenfunctions $`D^\theta `$ would have more singular form.
The above computations can be generalized to $`p`$-dimensions. Written in the complex coordinates $`z_i,\overline{z}_i`$ corresponding to oscillator operators (113), which diagonalize the noncommutativity matrix this looks as follows,
$$D_{\stackrel{}{m}\stackrel{}{n}}=D_{m_1n_1}^{\theta _{(1)}}(z_1,\overline{z}_1)D_{m_2n_2}^{\theta _{(2)}}(z_2,\overline{z}_2)\mathrm{}D_{m_{p/2}n_{p/2}}^{\theta _{(p/2)}}(z_{p/2},\overline{z}_{p/2}),$$
(117)
where,
$$[z_i,\overline{z}_j]_{}=\delta _{ij},i=1,\mathrm{},p/2.$$
(118)
Having the above map one can establish the following relations between operators and their Weyl symbols.
1. It is not difficult to derive that,
$$(2\pi )^{p/2}\sqrt{det\theta }tr๐
=dxF(x).$$
(119)
2. The (noncommutative) product of operators is mapped into the *star* or *Moyal* product of functions,
$$๐
๐FG(x),$$
(120)
where $`FG(x)`$ is defined as,
$$FG(x)=\mathrm{e}^{\frac{\mathrm{i}}{2}\theta ^{\mu \nu }_\mu _\nu ^{}}F(x)G(x^{})|_{x^{}=x}.$$
(121)
In terms of operator $`\widehat{\delta }(\widehat{x}x)`$, this product can be written as follows,
$$FG(x)=\mathrm{d}^py\mathrm{d}^pzK(x;y,z)F(y)G(z),$$
(122)
where,
$$\begin{array}{c}K(x;y,z)=\hfill \\ \hfill (2\pi )^{p/2}\sqrt{det\theta }tr\widehat{\delta }(\widehat{x}x)\widehat{\delta }(\widehat{x}y)\widehat{\delta }(\widehat{x}z)=\\ \hfill \mathrm{e}^{\frac{\mathrm{i}}{2}_\mu ^y\theta ^{\mu \nu }_\nu ^z}\delta (yx)\delta (zx),\end{array}$$
(123)
$`_\mu ^y`$ and $`_\mu ^z`$ are, respectively, $`/y^\mu `$ and $`/z^\mu `$, and in the last line one has ordinary delta functions.
On the other hand the ordinary product of functions was not found to have any reasonable meaning in this context.
3. One property of the star product is that in the integrand one can drop it once because of,
$$\mathrm{d}^pxFG(x)=\mathrm{d}^pxF(x)G(x),$$
(124)
were in the r.h.s the ordinary product is assumed.
4. Interesting feature of this representation is that partial derivatives of Weyl symbols correspond to commutators of respective operators with $`\mathrm{i}๐ฉ_\mu `$,
$$[\mathrm{i}๐ฉ_\mu ,๐
]\mathrm{i}(p_\mu FFp_\mu )(x)=\frac{F(x)}{x^\mu },$$
(125)
where $`p_\mu `$ is linear function of $`x^\mu `$: $`p_\mu =\theta _{\mu \nu }^1x^\nu `$.
This is an important feature of the star algebra of functions distinguishing it from the ordinary product algebra. In the last one can not represent the derivative as an *internal automorphism* while in the star algebra it is possible due to its nonlocal character. This property is of great importance in the field theory since, as it will appear later, it is the source of duality relations in noncommutative gauge models which we turn to in the next section.
###### Exercise 10
Derive equations (119)โ(125).
Let us turn back to the matrix model action (32) and represent an arbitrary matrix configuration as a perturbation of the background (69):
$$X_a=p_a+A_a,X_I=\mathrm{\Phi }_I,a=1,\mathrm{},p+1,I=p+2,\mathrm{},D.$$
(126)
Passing from operators $`A_a`$ and $`\mathrm{\Phi }`$ to their Weyl symbols using (108), (120) and (125) one gets following representation for the matrix action (32):
$$S=\mathrm{d}^px\left(\frac{1}{4}(_{ab}B_{ab})^2+\frac{1}{2}(_a\mathrm{\Phi }_I)^2\frac{1}{4}[\mathrm{\Phi }_I,\mathrm{\Phi }_J]_{}^2\right),$$
(127)
where,
$$F_{\mu \nu }=_\mu A_\nu _\nu A_\mu \mathrm{i}[A_\mu ,A_\nu ]_{}.$$
(128)
In the case of the irreducible representation of the algebra (70) this describes the U(1) gauge model.
One can consider a $`n`$-tuple degenerate representation in this case as well. As in the previous case the index labelling the representations become an internal symmetry index and the global gauge group of the model becomes U$`(n)`$. Indeed, the operator basis in which one can expand an arbitrary operator now is given by,
$$E_k^\alpha =\sigma ^\alpha \mathrm{e}^{\mathrm{i}k\widehat{x}},$$
(129)
where $`\sigma ^\alpha `$, $`\alpha =1,\mathrm{},n^2`$ are the adjoint generators of the u$`(n)`$ algebra. The can be normalized to satisfy,
$$[\sigma ^\alpha ,\sigma ^\beta ]=\mathrm{i}ฯต^{\alpha \beta \gamma }\sigma ^\gamma ,tr_{su(2)}\sigma ^\alpha \sigma ^\beta =\delta ^{\alpha \beta },$$
(130)
where $`ฯต^{\alpha \beta \gamma }`$ are the structure constants of the u$`(n)`$ algebra:
$$ฯต^{\alpha \beta \gamma }=\mathrm{i}tr_{su(2)}[\sigma ^\alpha ,\sigma ^\beta ]\sigma ^\gamma ,$$
(131)
which follows from (130). Then an operator $`\widehat{F}`$ is mapped to the following function $`F(x)`$:
$$\begin{array}{c}F^\alpha (x)=\sqrt{det\theta }\frac{\mathrm{d}^pk}{(2\pi )^{p/2}}\mathrm{e}^{\mathrm{i}kx}tr\left\{(\sigma ^\alpha \mathrm{e}^{\mathrm{i}k\widehat{x}})\widehat{F}\right\}=\hfill \\ \hfill (2\pi )^{p/2}\sqrt{det\theta }tr\left\{(\sigma ^\alpha \widehat{\delta }(\widehat{x}x))\widehat{F}\right\}.\end{array}$$
(132)
The equation (132) gives the most generic map from the space of operators to the space of $`p`$-dimensional u$`(n)`$-algebra valued functions.
###### Exercise 11
Prove that $`p`$ is always even.
Just for the sake of completeness let us give also the formula for the inverse map,
$$\widehat{F}=\mathrm{d}^px(\sigma ^\alpha \widehat{\delta }(\widehat{x}x))F^\alpha (x),$$
(133)
Applying the map (132) and (133) to the IKKT matrix model (32) or to the BFSS one (55), one gets, respectively, the $`p`$ or $`p+1`$ dimensional noncommutative u$`(n)`$ YangโMills model.
###### Exercise 12
Derive the $`p`$\- and $`(p+1)`$-dimensional noncommutative supersymmetric gauge model from the matrix actions (32) and (55), using the map (132) and its inverse (133).
Some comments regarding both gauge models described by the actions (87) and (128) are in order. In spite of the fact that both models look very similar to the โordinaryโ YangโMills models, the perturbation theory of this models are badly defined in the case of noncompact noncommutative spaces. In the first case the non-renormalizable divergence is due to extra integrations over $`l`$ in the โinternalโ space. In the case of noncommutative gauge model the behavior of the perturbative expansion is altered by the IR/UV mixing Minwalla:1999px ; VanRaamsdonk:2000rr . The supersymmetry or low dimensionality improves the situation allowing the โbadโ terms to cancel (see Sarkar:2002pb ; Bietenholz:2002ch ; Slavnov:2003ae ; Buric:2003qv ). On the other hand the compact noncommutative spaces provide both IR and UV cut off and the field theory on such spaces is finite Sheikh-Jabbari:1999iw . In the case of zero commutator background the behavior of the perturbative expansion depends on the eigenvalue distribution. Faster the eigenvalues increase, better the expansion converge. However there is always the problem of the zero modes corresponding to the diagonal matrix excitations (functions of commutative $`p_a`$โs). There is a hope that integrating over the remaining modes helps to generate a dynamical term for the zero modes too. Indeed, for purely bosonic model one has a repelling potential after the one-loop integration of the non-diagonal modes. The fermions contribute with the attractive potential. In the supersymmetric case the repelling bosonic contribution is cancelled by the attractive fermionic one and diagonal modes remain non-dynamical Makeenko:privat .
###### Exercise 13
Consider the EguchiโKawai model given by the action (9). Write down the equations of motion and find the classical solutions analogous to (69). One can have noncommutative solutions even for finite $`N`$. Explain, why? Consider arbitrary matrix configuration as a perturbation of the above classical backgrounds and find the resulting models. What is the space on which this models live? How the same space can be obtained from a non-compact matrix model.
We considered exclusively the bosonic models. When the supersymmetric theories are analyzed one has to deal also with the fermionic part. In the case of compact noncommutative spaces which correspond to finite size matrices one has a discrete system with fermions. In the lattice gauge theories with fermions there is a famous problem related to the fermion *doubling* Nielsen:1981hk . Concerning the theories on the compact noncommutative spaces it was found that in some cases one can indeed have fermion doubling Sochichiu:2000fs <sup>4</sup><sup>4</sup>4For the case of the unitary EguchiโKawai-type model with fermions see Kitsunezaki:1997iu some other cases were reported to be doubling free and giving alternative solutions to the long standing lattice problem Balachandran:2003ay .
## 6 Matrix models and dualities of noncommutative gauge models
In the previous section we realized that the matrix model from different โpointsโ of the moduli space of classical solutions looks like different gauge models. These models can have different dimensionality or different global gauge symmetry group, but they all are equivalent to the original IKKT or BFSS matrix model. This equivalence can be used to pass from some noncommutative model back to the matrix model and then to a different noncommutative model and viceversa. Thus, one can find a one-to-one map from one model to an equivalent one.
In reality one can jump the intermediate step by writing a new solution direct in the noncommutative gauge model and passing to Weyl (re)ordered description with respect to the new background. From the point of view of noncommutative geometry this procedure is nothing else that the change of the noncommutative variable taking into account also the ordering. Let us go to the details. Consider two different background solutions given by $`p_{\mu _{(i)}}^{(i)}`$, where $`\mu _{(i)}=1,\mathrm{},p_{(i)}`$ and the index $`i=1,2`$ labels the backgrounds. Denote the orders of degeneracy of the backgrounds by $`n_{(i)}`$. The commutator for both backgrounds is given by,
$$[p_{\mu _{(i)}}^{(i)},p_{\nu _{(i)}}^{(i)}]=\mathrm{i}B_{\mu _{(i)}\nu _{(i)}}^{(i)}.$$
(134)
Applying to a $`p_{(1)}`$-dimensional u$`(n_{(1)}`$ algebra valued field $`F^{\alpha _{(1)}}(x_{(1)})`$ first the inverse Weyl transformation (133) which maps it in the operator form and then the direct transformation (132) from the operator form to the second background one gets a $`p_{(2)}`$-dimensional u$`(n_{(2)}`$ algebra valued field $`F^{\alpha _{(2)}}(x_{(2)})`$ defined by
$$F^{\alpha _{(2)}}(x_{(2)})=\mathrm{d}^{p_{(1)}}x_{(1)}K_{(2|1)}^{\alpha _{(2)}\alpha _{(1)}}(x_{(2)}|x_{(1)})F^{\alpha _{(1)}}(x_{(1)}),$$
(135)
where the kernel $`K_{(2|1)}^{\alpha _{(2)}\alpha _{(1)}}(x_{(2)}|x_{(1)})`$ is given by,
$$\begin{array}{c}K_{(1|2)}^{\alpha _{(2)}\alpha _{(1)}}(x_{(2)},x_{(1)})=(2\pi )^{p_{(2)}/2}\sqrt{det\theta _{(2)}}\times \hfill \\ \hfill tr\left\{(\sigma _{(2)}^{\alpha _{(2)}}\widehat{\delta }(\widehat{x}_{(2)}x_{(2)}))(\sigma _{(1)}^{\alpha _{(1)}}\widehat{\delta }(\widehat{x}_{(1)}x_{(1)}))\right\},\end{array}$$
(136)
where $`x_{(i)}`$ and $`\sigma _{(i)}^{\alpha _{(i)}}`$ are the coordinate and algebra generators corresponding to the background $`p_{\mu _{(i)}}^{(i)}`$.
The equation (136) still appeals to the background independent operator form by using the $`\widehat{\delta }`$-operators and trace. This can be eliminated in the following way. Consider the functions $`x_{(2)}^{\mu _{(2)}}(x_{(1)}^{\mu _{(1)}},\sigma _{(1)}^{\alpha _{(1)}})=x_{(2)}^{\mu _{(2)};\alpha _{(1)}}(x_{(1)}^{\mu _{(1)}})\sigma _{(1)}^{\alpha _{(1)}}`$ and $`\sigma _{(2)}^{\alpha _{(2)}}(x_{(1)}^{\mu _{(1)}},\sigma _{(1)}^{\alpha _{(1)}})=\sigma _{(2)}^{\alpha _{(2)};\alpha _{(1)}}(x_{(1)}^{\mu _{(1)}})\sigma _{(1)}^{\alpha _{(1)}}`$ which are the symbols of the second background $`\widehat{x}_{(2)}^{\mu _{(2)}}`$ which are Weyl-ordered with respect to the first background. Namely, they are the solution to the equation,
$$x_{(2)}^{\mu _{(2)}}_{(1)}x_{(2)}^{\nu _{(2)}}x_{(2)}^{\nu _{(2)}}_{(1)}x_{(2)}^{\mu _{(2)}}=\theta _{\mu _{(2)}\nu _{(2)}}^{(2)},$$
(137)
and for $`\sigma _{(2)}`$
$$\sigma _{(2)}^{\alpha _{(2)}}_{(1)}\sigma _{(2)}^{\beta _{(2)}}\sigma _{(2)}^{\alpha _{(2)}}_{(1)}\sigma _{(2)}^{\beta _{(2)}}=\mathrm{i}ฯต^{\alpha _{(2)}\beta _{(2)}\gamma _{(2)}}\sigma _{(2)}^{\gamma _{(2)}}$$
(138)
where $`_{(1)}`$ includes both the noncommutative with $`\theta _{(1)}`$ and the u$`(n_{(1)})`$ matrix products and we did not write explicitly the arguments $`(x_{(1)}^{\mu _{(1)}},\sigma _{(1)}^{\alpha _{(1)}})`$ and u$`(n_{(1)})`$ matrix indices of $`x_{(2)}`$ and $`\sigma _{(2)}`$. Then, the kernel (136) can be rewritten in the $`x_{(1)}`$ background as follows,
$$\begin{array}{c}K_{(1|2)}^{\alpha _{(2)}\alpha _{(1)}}(x_{(2)},x_{(1)})=\hfill \\ \hfill \sqrt{\frac{det2\pi \theta _{(2)}}{det2\pi \theta _{(1)}}}d_{(1)}^{\alpha _{(1)}\beta _{(1)}\gamma _{(1)}}\left(\sigma _{(2)}^{\alpha _{(2)};\beta _{(1)}}_{(1)}\delta _{_{(1)}}^{\gamma _{(1)}}(x_{(2)}(x_{(1)})x_{(2)})\right),\end{array}$$
(139)
where $`d_{(1)}^{\alpha \beta \gamma }=tr_{(1)}\sigma _{(1)}^\alpha \sigma _{(1)}^\beta \sigma _{(1)}^\gamma `$ and
$$\delta _{_{(1)}}^{\gamma _{(1)}}(x_{(2)}(x_{(1)})x_{(2)})=\frac{\mathrm{d}^{p_{(2)}}l}{(2\pi )^{p_{(2)}}}tr_{(1)}\sigma _{(1)}^{\gamma _{(1)}}\mathrm{e}_{_{(1)}}^{\mathrm{i}l(x_{(2)}(x_{(1)})x_{(2)})},$$
(140)
$`\mathrm{e}_{}^{f(x)}`$ is the star exponent computed with the noncommutative structure corresponding to $``$.
General expression for the basis transform (135) with the kernel (136) or (139) looks rather complicate almost impossible to deal with. Therefore it is useful to consider some particular examples which we take from Kiritsis:1997hj which show that in fact the objects are still treatable.
### 6.1 Example 1: The U(1) $``$ U($`n`$) map
Let us present the explicit construction for the map from U(1) to U(2) gauge model in the case of two-dimensional non-commutative space. The map we are going to discuss can be straightforwardly generalised to the case of arbitrary even dimensions as well as to the case of arbitrary U(n) group.
The two-dimensional non-commutative coordinates are,
$$[x^1,x^2]=\mathrm{i}\theta .$$
(141)
As we already discussed, non-commutative analog of complex coordinates is given by oscillator rising and lowering operators,
$$a=\sqrt{\frac{1}{2\theta }}(x^1+\mathrm{i}x^2),\overline{a}=\sqrt{\frac{1}{2\theta }}(x^1\mathrm{i}x^2)$$
(142)
$$a|n=\sqrt{n}|n1,\overline{a}|n=\sqrt{n+1}|n+1,$$
(143)
where $`|n`$ is the oscillator basis formed by eigenvectors of $`N=\overline{a}a`$,
$$N|n=n|n.$$
(144)
The gauge symmetry in this background is non-commutative U(1).
We will now construct the non-commutative U(2) gauge model. For this, consider the U(2) basis which is given by following vectors,
$$|n^{},a=|n^{}e_a,a=0,1$$
(145)
$$e_0=\left(\begin{array}{c}1\\ 0\end{array}\right),e_1=\left(\begin{array}{c}0\\ 1\end{array}\right),$$
(146)
where $`\{|n^{}\}`$ is the oscillator basis and $`\{e_a\}`$ is the โisotopicโ space basis.
The one-to-one correspondence between U(1) and U(2) bases can be established in the following way Nair:2001rt ,
$$|n^{}e_a|n=|2n^{}+a,$$
(147)
where $`|n`$ is a basis element of the U(1)-Hilbert space and $`|n^{}e_a`$ is a basis element of the Hilbert space of U(2)-theory. (Note, that they are two bases of the same Hilbert space.)
Let us note that the identification (147) is not unique. For example, one can put an arbitrary unitary matrix in front of $`|n`$ in the r.h.s. of (147). This in fact describes all possible identifications and respectively maps from U(1) to U(2) model.
Under this map, the U(2) valued functions can be represented as scalar functions in U(1) theory. For example, constant U(2) matrices are mapped to particular functions in U(1) space. To find these functions, it suffices to find the map of the basis of the u$`(2)`$ algebra given by Pauli matrices $`\sigma _\alpha `$, $`\alpha =0,1,2,3`$.
In the U$`(1)`$ basis Pauli matrices look as follows,
$`\sigma _0={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left(|2n2n|+|2n+12n+1|\right)๐,`$ (148a)
$`\sigma _1={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left(|2n2n+1|+|2n+12n|\right),`$ (148b)
$`\sigma _2=\mathrm{i}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left(|2n2n+1||2n+12n|\right),`$ (148c)
$`\sigma _3={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left(|2n2n||2n+12n+1|\right),`$ (148d)
while the โcomplexโ coordinates $`a^{}`$ and $`\overline{a}^{}`$ of the U(2) invariant space are given by the following,
$`a^{}={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\sqrt{n}\left(|2n22n|+|2n12n+1|\right),`$ (149a)
$`\overline{a}^{}={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\sqrt{n+1}\left(|2n+22n|+|2n+32n+1|\right).`$ (149b)
One can see that when trying to find the Weyl symbols for operators given by (148), (149), one faces the problem that the integrals defining the Weyl symbols diverge. This happens because the respective functions (operators) do not belong to the non-commutative analog of $`L^2`$ space (are not square-trace).
Let us give an alternative way to compute the functions corresponding to operators (148) and (149). To do this let us observe that operators
$$\mathrm{\Pi }_+=\underset{n=0}{\overset{\mathrm{}}{}}|2n2n|,$$
(150)
and
$$\mathrm{\Pi }_{}=\underset{n=0}{\overset{\mathrm{}}{}}|2n+12n+1|,$$
(151)
can be expressed as<sup>5</sup><sup>5</sup>5Weyl symbols of $`a`$ and $`\overline{a}`$ are denoted, respectively, as $`z`$ and $`\overline{z}`$. The same rule applies also to primed variables.
$$\mathrm{\Pi }_+=\frac{1}{2}\underset{n=0}{\overset{\mathrm{}}{}}\left(1+\mathrm{sin}\pi \left(n+\frac{1}{2}\right)\right)|nn|\frac{1}{2}\left(1+\mathrm{sin}_{}\pi \left(\overline{z}z+\frac{1}{2}\right)\right),$$
(152)
and,
$$\mathrm{\Pi }_{}=\mathrm{i}\mathrm{\Pi }_+=\frac{1}{2}\left(1\mathrm{sin}_{}\pi \left(\overline{z}z+\frac{1}{2}\right)\right)=\frac{1}{2}\left(1\mathrm{sin}_{}\pi |z|^2\right),$$
(153)
where $`\mathrm{sin}_{}`$ is the โstarโ sin function defined by the star Taylor series,
$$\mathrm{sin}_{}f=f\frac{1}{3!}fff+\frac{1}{5!}fffff\mathrm{},$$
(154)
with the star product defined in variables $`z,\overline{z}`$ as follows,
$$fg(\overline{a},a)=\mathrm{e}^{\overline{}^{}\overline{}^{}}f(\overline{z},z)g(\overline{z}^{},z^{})|_{z^{}=z},$$
(155)
where $`=/z`$, $`\overline{}=/\overline{z}`$ and analogously for primed $`z^{}`$ and $`\overline{z}^{}`$. For convenience we denoted Weyl symbols of $`a`$ and $`\overline{a}`$ as $`z`$ and $`\overline{z}`$.
The easiest way to compute (152) and (153) is to find the Weyl symbol of the operator,
$$I_k^\pm =\frac{1\pm \mathrm{sin}\left(\overline{a}a+\frac{1}{2}\right)}{(\overline{a}a+\gamma )^k},$$
(156)
were $`\gamma `$ is some constant, mainly $`\pm 1/2`$.
For sufficiently large $`k`$, the operator $`I_k^\pm `$ becomes square trace for which the formula (132) defining the Weyl map is applicable. The result can be analytically continued for smaller values of $`k`$, using the following recurrence relation,
$$I_{km}^\pm (\overline{z},z)=\underset{m\text{ times}}{\underset{}{\left(|z|^2+\gamma \frac{1}{2}\right)\mathrm{}\left(|z|^2+\gamma \frac{1}{2}\right)}}I_k^\pm (\overline{z},z).$$
(157)
The last equation requires computation of only finite number of derivatives of $`I_k^\pm (\overline{z},z)`$ arising from the star product with polynomials in $`\overline{z},z`$.
###### Exercise 14
Compute the Weyl symbol for the operator (156).
### 6.2 Example 2: Map between different dimensions
Consider the situation when the dimension is changed. This topic was considered in Sochichiu:2000bg ; Sochichiu:2000kz .
Consider the Hilbert space $``$ corresponding to the representation of the two-dimensional non-commutative algebra (141), and $``$ (which is in fact isomorphic to $``$) which corresponds to the four-dimensional non-commutative algebra generated by
$$[x^1,x^2]=\mathrm{i}\theta _{(1)},[x^3,x^4]=\mathrm{i}\theta _{(2)}.$$
(158)
In the last case non-commutative complex coordinates correspond to two sets of oscillator operators, $`a_1`$, $`a_2`$ and $`\overline{a}_1`$, $`\overline{a}_2`$, where,
$`a_1=\sqrt{{\displaystyle \frac{1}{2\theta _{(1)}}}}(x^1+\mathrm{i}x^2),`$ $`\overline{a}_1=\sqrt{{\displaystyle \frac{1}{2\theta _{(1)}}}}(x^1\mathrm{i}x^2)`$ (159a)
$`a_1|n_1=\sqrt{n_1}|n_11,`$ $`\overline{a}_1|n_1=\sqrt{n_1+1}|n_1+1,`$ (159b)
$`a_2=\sqrt{{\displaystyle \frac{1}{2\theta _{(2)}}}}(x^3+\mathrm{i}x^4),`$ $`\overline{a}_2=\sqrt{{\displaystyle \frac{1}{2\theta _{(2)}}}}(x^3\mathrm{i}x^4)`$ (159c)
$`a|n_2=\sqrt{n_2}|n_21,`$ $`\overline{a}_2|n_2=\sqrt{n_2+1}|n_2+1,`$ (159d)
and the basis elements of the โfour-dimensionalโ Hilbert space $``$ are $`|n_1,n_2=|n_1|n_2`$.
The isomorphic map $`\sigma :`$ is given by assigning a unique number $`n`$ to each element $`|n_1,n_2`$ and putting it into correspondence to $`|n`$. So, the problems is reduced to the construction of an isomorphic map from one-dimensional lattice of e.g. nonnegative integers into the two-dimensional quarter-infinite lattice. This can be done by consecutive enumeration of the two-dimensional lattice nodes starting from the angle $`(00)`$. The details of the construction can be found in Refs. Sochichiu:2000bg ; Sochichiu:2000kz .
As we discussed earlier, this map induces an isomorphic map of gauge and scalar fields from two to four dimensional non-commutative spaces.
This can be easily generalized to the case with arbitrary number of factors $`\mathrm{}`$ corresponding to $`p/2`$ โtwo-dimensionalโ non-commutative spaces. In this way, one obtains the isomorphism $`\sigma `$ which relates two-dimensional non-commutative function algebra with a $`p`$-dimensional one, for $`p`$ even.
## 7 Example 3: Change of $`\theta `$
So far, we have considered maps which relate algebras of non-commutative functions in different dimensions or at least taking values in different Lie algebras. Due to the fact that they change considerably the geometry, these maps could not be deformed smoothly into the identity map. In this section we consider a more restricted class of maps which do not change either dimensionality or the gauge group but only the non-commutativity parameter. Obviously, this can be smoothly deformed into identity map, therefore one may consider infinitesimal transformations.
The new non-commutativity parameter is given by the solution to the equations of motion. In this framework, the map is given by the change of the background solution $`p_\mu `$ by an infinitesimal amount: $`p_\mu +\delta p_\mu `$. Then, a solution with the constant field strength $`F_{\mu \nu }^{(\delta p)}`$ will change the non-commutativity parameter as follows,
$$\theta ^{\mu \nu }+\delta \theta ^{\mu \nu }(\theta _{\mu \nu }^1+\delta \theta _{\mu \nu }^1)^1=(\theta _{\mu \nu }^1+F_{\mu \nu })^1.$$
(160)
Note, that the above equation does not require $`\delta \theta `$ to be infinitesimal.
Since we are considering solutions to the gauge field equations of motion $`A_\mu =\delta p_\mu `$ one should fix the gauge for it. A convenient choice would be e.g. the Lorentz gauge, $`_\mu \delta p_\mu =0`$. Then, the solution with
$$A_\mu ^{(\delta p)}\delta p_\mu =(1/2)ฯต_{\mu \nu }\theta ^{\nu \alpha }p_\alpha $$
(161)
with antisymmetric $`ฯต_{\mu \nu }`$ has the constant field strength
$$F_{\mu \nu }^{(\delta p)}\delta \theta _{\mu \nu }^1=ฯต_{\mu \nu }+(1/4)ฯต_{\mu \alpha }\theta ^{\alpha \beta }ฯต_{\beta \nu }=ฯต_{\mu \nu }+O(ฯต^2).$$
(162)
This corresponds to the following variation of the non-commutativity parameter,
$$\delta \theta ^{\mu \nu }=\theta ^{\mu \alpha }ฯต_{\alpha \beta }\theta ^{\beta \nu }\frac{1}{4}\theta ^{\mu \alpha }ฯต_{\alpha \gamma }\theta ^{\gamma \rho }ฯต_{\rho \beta }\theta ^{\beta \nu }=\theta ^{\mu \alpha }\delta \theta _{\alpha \beta }^1\theta ^{\beta \nu }+O(ฯต^2).$$
(163)
Let us note that such kind of infinitesimal transformations were considered in a slightly different context in Ishikawa:2001mq .
Let us find how non-commutative functions are changed with respect to this transformation. In order to do this, let us consider how the Weyl symbol (132) transforms under the variation of background (161). For an arbitrary operator $`\varphi `$ after short calculation we have,
$$\delta \varphi (x)=\frac{1}{4}\delta \theta ^{\alpha \beta }(_\alpha \varphi p_\beta (x)+p_\beta _\alpha \varphi (x)).$$
(164)
In obtaining this equation we had to take into consideration the variation of $`p_\mu `$ as well as of the factor $`\sqrt{det\theta }`$ in the definition of the Weyl symbol (132).
By the construction, this variation satisfies the โstar-Leibnitz ruleโ,
$$\delta (\varphi \chi )(x)=\delta \varphi \chi (x)+\varphi \delta \chi (x)+\varphi (\delta )\chi (x),$$
(165)
where $`\delta \varphi (x)`$ and $`\delta \chi (x)`$ are defined according to (164) and variation of the star-product is given by,
$$\varphi (\delta )\chi (x)=\frac{1}{2}\delta \theta ^{\alpha \beta }_\alpha \varphi _\beta \chi (x).$$
(166)
The property (165) implies that $`\delta `$ provides an homomorphism (which is in fact an isomorphism) of star algebras of functions.
The above transformation (164) do not apply, however, to the gauge field $`A_\mu (x)`$ and gauge field strength $`F_{\mu \nu }(x)`$. This is the case because the respective fields do not correspond to invariant operators. Indeed, according to the definition $`A_\mu =X_\mu p_\mu `$, where $`X_\mu `$ is corresponds to such an operator. Therefore, the gauge field $`A_\mu (x)`$ transforms in a nonhomogeneous way,<sup>6</sup><sup>6</sup>6In fact the same happens in the map between different dimensions.
$$\delta A_\mu (x)=\frac{1}{4}\delta \theta ^{\alpha \beta }(_\alpha A_\mu p_\beta +p_\beta _\alpha A_\mu )+\frac{1}{2}\theta _{\mu \alpha }\delta \theta ^{\alpha \beta }p_\beta .$$
(167)
The transformation law for $`F_{\mu \nu }(x)`$ can be computed using its definition (128) and the โstar-Leibnitz ruleโ (165) as well as the fact that it is the Weyl symbol of the operator,
$$F_{\mu \nu }=\mathrm{i}[X_\mu ,X_\nu ]\theta _{\mu \nu }.$$
(168)
Of course, both approaches give the same result,
$$\delta F_{\mu \nu }(x)=\frac{1}{4}\delta \theta ^{\alpha \beta }(_\alpha F_{\mu \nu }p_\beta +p_\beta _\alpha F_{\mu \nu })(x)\delta \theta _{\mu \nu }^1.$$
(169)
The infinitesimal map described above has the following properties:
1. It maps gauge equivalent configurations to gauge equivalent ones, therefore it satisfies the SeibergโWitten equation,
$$U^1AU+U^1dUU^1^{}A^{}^{}U^{}+U^1^{}d^{}U^{}.$$
(170)
2. It is linear in the fields.
3. Any background independent functional is invariant under this transformation. In particular, any gauge invariant functional whose dependence on gauge fields enters through the combination $`X_{\mu \nu }(x)=F_{\mu \nu }+\theta _{\mu \nu }^1`$ is invariant with respect to (164)โ(169). This is also the symmetry of the action provided that the gauge coupling transforms accordingly.
4. Formally, the transformation (164) can be represented in the form,
$$\delta \varphi (x)=\delta x^\alpha _\alpha \varphi (x)=\varphi (x+\delta x)\varphi (x),$$
(171)
where $`\delta x^\alpha =\theta ^{\alpha \beta }\delta p_\beta `$ and no star product is assumed. This looks very similar to the coordinate transformations.
The map we just constructed looks very similar to the famous SeibergโWitten map, which is given by the following variation of the background $`p_\mu `$ Seiberg:1999vs ,
$$\delta _{\mathrm{SW}}p_\mu =\frac{1}{2}ฯต_{\mu \nu }\theta ^{\nu \alpha }A_\alpha .$$
(172)
In (161) we have chosen $`\delta p_\mu `$ independent of gauge field background. (In fact the gauge field background was switched-on later, after the transformation.) An alternative way would be to have nontrivial field $`A_\mu (x)`$ from the very beginning and to chose $`\delta p_\mu `$ to be of the SeibergโWitten form. Then, the transformation laws corresponding to such a transformation of the background coincide exactly with the standard SW map. This appears possible because the function $`p_\mu =\theta _{\mu \nu }^1x^\nu `$ has the same gauge transformation properties as $`A_\mu (x)`$,
$$p_\mu U^1p_\mu U(x)U^1_\mu U(x).$$
(173)
## 8 Discussion and outlook
This lecture notes were designed as a very basic and very subjective introduction to the field. Many important things were not reflected and even not mentioned here. Among these, very few was said about the brane dynamics and interpretation which was the main motivation for the development of the matrix models, while the literature on this topic is enormously vast. For this we refer the reader to other reviews and lecture notes mentioned in the introduction (as well as to the references one can find inside these papers).
Recently, the role of the matrix models in the context of AdS/CFT correspondence became more clear. Some new matrix models arise in the description of the anomalous dimensions of composite super-YangโMills operators (see e.g Agarwal:2004cb ; Bellucci:2004fh .
Another recent progress even not mentioned here but which is related to matrix models is their use for the computation of the superpotential of $`๐ฉ=1`$ supersymmetric gauge theories Dijkgraaf:2002dh ; Dijkgraaf:2002vw ; Dijkgraaf:2002fc .
Acknowledgements. This lecture notes reflects the experience I gained due to the communication with many persons. My thanks are directed most of all to my collaborators and colleagues from Bogoliubov lab in Dubna, Physics Dept of University of Crete, Laboratori Nazionali di Frascati. The complete list of persons is too long.
I am grateful to the friends who helped me with various problems a while this text was being written: Giorgio Pagnini, Carlo Cavallo, Valeria and Claudio Minardi.
This work was supported by the INTAS-00-00262 grant and the Alexander von Humboldt research fellowship.
## Index
* $`1`$-brane ยง2.2
* โthe second string revolutionโ ยง2.1
* $`0`$-branes ยง2.1
* area preserving diffeomorphisms ยง4.1
* BFSS matrix model ยง1, ยง2.3, ยง3.2, ยง3.2
* BMN matrix model ยง2.3
* brane ยง1, ยง2.1
* ChanโPatton factor ยง2.1
* chemical potential ยง2.2
* compact branes ยง4.2
* D-string tension ยง3.1
* EguchiโKawai model ยง2.2
* fermionic part ยง3, ยง3.2
* fuzzy sphere ยง3.1
* IKKT matrix model ยง1, ยง2.2, ยง3.1
* M-theory ยง1
* mass term ยง4.1
* matrix mechanics ยง1
* membrane ยง3.2
* $`N`$ ยง2.2
* NambuโGoto action ยง3.1, ยง3.2
* NambuโGotoโPolyakov equations of motion ยง4.1
* noncommutative torus ยง3.1, ยง3.1
* oscillator basis ยง3.1
* Poisson bracket ยง3.1, ยง3.2
* Polyakovโs trick ยง3.1, ยง3.2
* random matrices ยง1, ยง2.2
* reparametrization invariance ยง3.2
* spherical branes ยง2.3
* star product ยง3.1
* Weyl symbol ยง3.1
* worldsheet deformation ยง3
* worldsheet quantization ยง3.1
* worldvolume deformation ยง3
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# 1 Introduction
## 1 Introduction
### 1.1 Motivation
Which one is the exception- a motor boat, a car, an airplane or a submarine? The immediate (perhaps, after a short contemplation) answer is, of course, the car. Indeed, a car is the only vehicle in the above list which โdefiesโ the law of conservation of linear momentum. The others create an opposite stream in the medium in which they move (air, water) which, by conservation of linear momentum, must be compensated by the motion of the vehicle. Take away the friction of these vehicles and the medium, and they will perform better (faster, more efficient).
The principle behind the motion of a car is different. It is moving (that is, shifting from rest to cruise velocity) because of the friction created by the contact of its tires with the road. Take away the friction and the car will not be able to move at all, independently of how hard you push the gas pedal. You will not be able to stop either, if your friend, unwisely attempting to help, gave you an initial push on a frictionless road.
So, what is a car? If we strip it off the non-essential components (radio, GPS, the fashioned seat covers etc.), it is a collection of components which can move with respect to each other under the preassign constraints caused by the mechanical structure. Unless you change gear or the pressure you apply on the gas pedal, the motion of these inner part can be assumed to be periodic (again, with respect to the frame of reference of the car itself). This concept has a lot in common with the subject of optimal locomotion of a swimmer in a Stokes flow - where the motion is due to a periodic change of shape of the swimmer in the absence of inertia. The description below is, in a sense, the mechanical analog of the swimmer model, see \[AGK\] and the references listed there. Another series of publications which seems to be related to the present discussion concerns molecular motors and the flashing rachet. See \[CHK\], \[CKK\], DKK\] and ref. therein.
Note that in the case of micro-swimmers and molecular motors, it seems that there is no intuitive way to predict the direction and velocity of the swimmer (res. motor) from its periodic motion. A car, at a first glance, is different. However, the abstraction of a car which we consider in this paper (and call a โmechanical systemโ) is, in a way, a generalization of the concept which contains cars, microswimers, molecular motors and, perhaps, many other objects whose dynamics are not inertial in nature.
### 1.2 Objectives and outline of results
A mechanical system is represented by a Lagrangian $`L=L(\dot{x},t)`$. It is $`T`$periodic in the time $`t`$, and $`x(t)`$ stands for the position of a reference point (say, the center of mass of the system). Such a system is called โmobilizedโ if the global minimizer $`x(t)`$ of the action $`_0^TL(\dot{x},t)๐t`$ is not periodic, i.e $`x(T)x(0)`$. In section 2 we attempt to justify this model and explain the reason why the global minimum of the action represents the asymptotic motion of the system under friction. In addition, we introduce a reasonable definition of efficiency, in terms of a relation involving the speed $`\overline{v}:=|x(T)x(0)|/T`$ due to the action minimizer, and the minimal action $`\overline{D}=_0^TL(\dot{x},t)๐t`$ itself. The number $`\overline{D}`$ stands for some indication of the energy dissipated (or invested) per period, so a more efficient system means larger $`\overline{v}`$ and smaller $`\overline{D}`$. We scale the efficiency function $`e_L`$ of a mechanical system $`L`$ in such a way that $`0<e_L<1`$. Then, we ask if either there is a (theoretical) possibility to achieve the โmost efficientโ system $`\overline{L}`$ whose efficiency $`\overline{e}:=e_{\overline{L}}<1`$, or there is a sequence of systems $`L_n`$ whose efficiencies $`e_{L_n}\overline{e}1`$, but the โidealโ, most efficient system $`\overline{L}`$ does not exist.
In section 3 we consider a special case of mechanical systems, where $`L`$ is a homogeneous function of $`\dot{x}`$, and prove the first alternative: There is a โbestโ (most efficient) mechanical system and its efficiency $`\overline{e}`$ is always smaller than $`1`$.
## 2 Description of the model
Let us attempt to build a mathematical caricature the โcarโ, composed of a finite number of โpartsโ. To wit, assume it is composed of a collection of $`N`$ points of respective masses $`m_i`$ executing orbits $`x_i=x_i(t)`$ on the line $``$, $`i=1,\mathrm{}n`$, so that
$$x_i(t+T)=x_i(t)+T\overline{v}$$
(2.1)
where $`\overline{v}`$ is the effective velocity of the car (with respect to a reference frame attached to the road) and $`T`$ is the period of one cycle.
We assume $`_1^Nm_i=1`$. The orbit $`x_i(t)`$ is given with respect to a fixed frame of reference (attached to the road). We may write it as
$$x_i(t)=y_i(t)+x(t)$$
where $`y_i`$ is the orbit of the corresponding point with respect to a reference attached to the car and $`x=x(t)`$ is the orbit of the car as a whole with respect to a reference attached to the road. In this representation $`y_i`$ is a periodic orbit, representing the motion of a part forced by, say, the internal combustion of the engine.
The simplest model for the motion of $`x(t)`$ is the linear forced system
$$\ddot{x}+\beta \dot{x}=F(t)$$
(2.2)
where $`\beta >0`$ is the friction coefficient and $`F`$ is the total forcing acted on the car by the motion of its inner parts. It is the sum of the forces $`F:=_i^Nf_i`$, where $`f_i`$, the force applied by the $`i`$ part, is defined in terms of $`y_i`$ as:
$$f_i(t)=m_i(\ddot{y}_i+\beta \dot{y}_i).$$
(2.3)
Now, define
$$^{\{y\}}(\dot{x},t):=\underset{i=1}{\overset{N}{}}m_iL(\dot{y}_i(t)+\dot{x}).$$
(2.4)
Then (2.2, 2.3) can be summarized as the Euler-Lagrange equation corresponding to the Lagrangian
$$_\beta ^{\{y\}}(\dot{x},t):=e^{\beta t}^{\{y\}}(\dot{x},t)$$
(2.5)
where
$$L(s):=|s|^2/2.$$
(2.6)
Indeed, the Euler-Lagrange equation associated with $`_\beta ^{\{y\}}`$ under (2.6) is
$$0=e^{\beta t}\left[\underset{1}{\overset{N}{}}m_i(\ddot{x}+\ddot{y}_i)+\beta \underset{1}{\overset{N}{}}m_i(\dot{x}+\dot{y}_i)\right]$$
which implies (2.2) via (2.3) and the condition $`m_i=1`$.
The energy dissipated per cycle for an orbit $`y_1,\mathrm{}y_N,x`$ is given by $`\beta D`$ where
$$D:=\underset{i=1}{\overset{N}{}}m_i_0^T\frac{|\dot{y}_i+\dot{x}|^2}{2}๐t=_0^T^{\{y\}}(\dot{x},t)๐t.$$
(2.7)
We now generalize (2.4-2.7) into:
###### Definition 2.1.
A mechanical system $`๐_\beta `$ is determined by a forced orbit composed of $`N`$ periodic functions $`๐ฒ=\{y_i(t),\mathrm{}y_N(t)\}`$ in terms of the Lagrangian $`_\beta ^{\{y\}}`$ as given in (2.5) and a convex function $`L`$ generalizing (2.6).
The moment associated with $`^{\{y\}}`$ is denoted by $`p=_{\dot{x}}^{\{y\}}`$. The Euler-Lagrange equation associated with (2.5) is
$$\dot{p}+\beta p=0p(t)0$$
so the asymptotic motion of a mechanical system $`๐_\beta `$ is determined by the orbit $`p(t)0`$.
To elaborate, let
$$^{\{y\}}(p,t)=\underset{\zeta }{sup}\left[p\zeta ^{\{y\}}(\zeta ,t)\right].$$
(2.8)
be the Hamiltonian associated with the Lagrangian $`^{\{y\}}`$. The equation of motion corresponding to the non-dissipative dynamics is given by
$$\dot{x}(t)=_p^{\{y\}}(\lambda ,t)$$
where $`\lambda `$ stands for the constant momentum $`p`$. If a friction is applied, then $`x=x(t)`$ is determined by $`p=0`$, namely
$$\dot{x}(t)=_p^{\{y\}}(0,t)$$
is the asymptotic motion of the system. It is a global minimizer of the action determined by the Lagrangian $`^{\{y\}}`$. Note that
$$\underset{x=x(t)}{\mathrm{min}}_0^T^{\{y\}}(\dot{x},t)๐t=_0^T^{\{y\}}(0,t)๐t.$$
Let
$$\overline{v}(0)=\frac{1}{T}_0^T_p^{\{y\}}(0,t)๐t\frac{x(T)x(0)}{T},$$
where $`x(t)`$ is the global minimizer of the action.
###### Definition 2.2.
A mechanical system for which $`\overline{v}(0)0`$ is called a mobilized system.
The first result is somewhat disappointing:
###### Theorem 1.
If $`L`$ is a quadratic function (2.6), then the system is not mobilized .
###### Proof.
$$^{\{y\}}(\dot{x},t)=\frac{1}{2}\underset{i=1}{\overset{N}{}}m_i\left|\dot{y}_i(t)+\dot{x}\right|^2_{\dot{x}}^{\{y\}}=\underset{i=1}{\overset{N}{}}m_i\left(\dot{y}_i(t)+\dot{x}\right).$$
In particular, $`_{\dot{x}}^{\{y\}}=0`$ implies $`\dot{x}=_{i=1}^Nm_i\dot{y}_i`$. Since $`y_i`$ are periodic by definition, then $`T^1_0^T\dot{x}๐t=\overline{v}(0)=0`$. โ
We now generalize the energy dissipation in the linear case (2.7) for general $`๐_\beta `$ mechanical systems. We shall denote by $`\overline{D}`$ the minimal action, and refer to it as the energy dissipated along a cycle:
$$\overline{D}=\underset{x=x(t)}{\mathrm{min}}\frac{1}{T}_0^T^{\{y\}}(\dot{x},t)๐t=\frac{1}{T}_0^T^{\{y\}}(0,t)๐t.$$
(2.9)
Next, we define the efficiency of a mobilized system $`๐:=y_1,\mathrm{}y_N`$. This should indicate the ratio of the distance transversed per cycle to the dissipated energy. The right scaling for this turns out to be
$$e_L(y_1,\mathrm{}y_N):=\frac{L\left(\overline{v}(0)\right)}{\overline{D}}=\frac{L\left(\frac{1}{T}_0^T_p^{\{y\}}(0,t)๐t\right)}{T^1_0^T^{\{y\}}(0,t)๐t}.$$
(2.10)
In fact, it can be proven that
###### Lemma 2.1.
If $`L`$ is a convex function, then for any mechanical system composed of $`N`$ periodic orbits $`y_i(t)=y_i(T+t),i=1,\mathrm{}N`$ ,
$$0<e_L(y_1,\mathrm{}y_N)1.$$
If, moreover, $`L`$ is strictly convex, then $`e_L<1`$.
###### Proof.
Let $`x(t)`$ be the any orbit. Then,
$$T^1_0^T^{\{y\}}(\dot{x},t)๐t=T^1\underset{i=1}{\overset{N}{}}m_i_0^TL(\dot{y}_i(t)+\dot{x}(t))๐t.$$
By Jensenโs inequality, the normalization condition $`m_i=1`$, the periodicity of $`y_i`$ and the convexity of $`L`$:
$$T^1\underset{i=1}{\overset{N}{}}m_i_0^TL(\dot{y}_i(t)+\dot{x}(t))L\left(T^1_0^T\underset{i}{\overset{N}{}}m_i(\dot{y}_i+\dot{x})dt\right)=L\left(T^1_0^T\dot{x}๐t\right),$$
(2.11)
so $`0<e_L1`$. Now, if $`L`$ is strictly convex then the equality in (2.11) holds if and only if $`N=1`$ or $`y_iy_j`$ for all $`1i,jN`$. But, in the later cases, the optimal orbit $`x`$ is evidently equal to any component $`y_i`$, so it is a periodic function and $`\overline{v}(0)=0`$. Hence the system is not mobilized and efficiency is not defined (or $`e_L=0`$ altogether). โ
Assume now that a convex Lagrangian $`L`$ is given, as well as $`N`$ and $`\{m_1,\mathrm{}m_N\}^{+,N}`$, $`_1^Nm_i=1`$. Let
$$\mathrm{\Lambda }_L:=\{๐=(y_1,\mathrm{}y_N)C^1([0,1];^N);๐(0)=๐(1).\}$$
(2.12)
Let
$$\overline{e}_L:=\underset{๐\mathrm{\Lambda }}{sup}e_L(๐).$$
(2.13)
We know that $`\overline{e}_L1`$. The intriguing questions are
i) Is $`\overline{e}_L=1`$?
ii) If $`\overline{e}_L<1`$, can the supremum in (2.13) be achieved in some sense?
We try to answer these questions in the special case of homogeneous Lagrangians.
## 3 Homogeneous Lagrangians
To fix the idea, let us concentrate on the case $`L(s)=|s|^\sigma `$ for some $`\sigma >1`$. We shall further assume a unit period $`T=1`$.
Let
$$H^\sigma :=\{x=x(t):[0,1],0t1;_0^1\left|\dot{x}\right|^\sigma <\mathrm{}\}.$$
(3.1)
Let also
$$\mathrm{\Lambda }_\sigma :=\left\{๐\left(H^\sigma \right)^N,๐(0)=๐(1)\right\}.$$
(3.2)
Given $`m_i=m(0,1)`$, $`_1^Nm_i=1`$ and $`๐=(y_1,\mathrm{}y_N)\mathrm{\Lambda }_\sigma `$ consider the Lagrangian
$$(\dot{x},\dot{๐}):=\underset{i=1}{\overset{N}{}}m_i|\dot{y}_i+\dot{x}|^\sigma .$$
(3.3)
Define also for any such $`๐`$ and $`p^N`$,
$$(p,\dot{๐}):=\underset{\xi ^N}{sup}\left\{p\xi (\xi ,\dot{๐})\right\},$$
(3.4)
and
$$๐^๐(p):=_0^1(p,\dot{๐}(t))๐t$$
(3.5)
Also, define:
$$u(\dot{๐}):=\frac{(p,\dot{๐})}{p}|_{p=0}$$
(3.6)
and
$$<u>_๐:=_0^1u(\dot{๐}(t))๐t.$$
(3.7)
We obtain via (3.4-3.7), (2.10) and (2.13)
$$e_\sigma (๐):=\frac{\left|<u>_๐\right|^\sigma }{๐^๐(0)};\overline{e}_\sigma =\underset{๐\mathrm{\Lambda }_\sigma }{sup}e_\sigma (๐).$$
(3.8)
One may wonder if, under the homogeneity condition, the mechanical system is mobilized, i.e $`\overline{e}_\sigma >0`$. The first result we claim is that this condition is equivalent to a property of the function $`u`$ as defined in (3.6), namely
###### Lemma 3.1.
The homogeneous system is mobilized if and only if $`u`$ is not a linear function.
###### Proof.
If $`u`$ is a linear function, then for any $`๐\mathrm{\Lambda }_\sigma `$, $`u\left(\dot{๐}\right)๐t=0`$ by definition of $`\mathrm{\Lambda }_\sigma `$. Indeed, the condition $`๐(0)=๐(1)`$ is equivalent to $`_0^1\dot{๐}(t)๐t=0`$.
Conversely, if $`u`$ is not a linear function, then there exists $`๐\mathrm{\Lambda }_\sigma `$ for which $`_0^1u\left(\dot{๐}\right)๐t0`$. This implies $`e_\sigma (๐)>0`$ and, in particular, the system is mobilized. โ
The condition of non-linearity of $`u`$ is rather delicate. In fact,
###### Lemma 3.2.
If either $`\sigma =2`$ or $`N=2`$ then $`u`$ is a linear function.
###### Proof.
We already know that $`\sigma =2`$ (quadratic Lagrangian) is not mobilized for any $`N`$ by Theorem 1.
By the homogeneity of $`L`$ and the definition of $`u`$ we observe that $`u`$ satisfies the pair of symmetries:
$$\zeta ^N,๐=(1,\mathrm{}1),and\lambda ,u(๐ป+\lambda ๐)=u(๐ป)+\lambda $$
$$\alpha ^+,u(\alpha ๐ป)=\alpha u(๐ป)$$
(3.9)
We conclude, therefore, that $`u(y_1,y_2)=f(y_1y_2)+y_1`$ holds for some function $`f`$ of a single variable. From the second equality of (3.9) it follows that
$$f(\alpha (y_1y_2))+\alpha y_1=\alpha f(y_1y_2)+\alpha y_1\alpha f(\zeta )=f(\alpha \zeta )$$
for any $`\alpha `$ and $`\zeta `$. Hence, $`f`$ is linear and so is $`u`$. โ
The main result of this paper is:
### 3.1 Main result
###### Theorem 2.
If the homogeneous mechanical system (3.3) is mobilized then there exists a maximizer of $`e_\sigma `$ in the set $`\mathrm{\Lambda }_\sigma `$, and $`\overline{e}_\sigma <1`$.
Remark: Lemma 3.2 implies that $`N>2`$ and $`\sigma 2`$ are necessary for the condition of Theorem 2. We conjecture that these conditions are also sufficient. In any case, it is not difficult to construct examples for mobilized homogeneous systems. For example, take $`\sigma =3`$, $`N=3`$ and $`m_1=m_2=m_3=1/3`$. The function $`u=u(y_1,y_2,y_3)`$ can be readily calculated as the root of a quadratic equation whose coefficients are linear functions of $`y_i`$. The discriminant, however, is not a complete square, so $`u`$ is not linear.
In the rest of this section we prove Theorem 2.
From (3.9) we obtain that $`(0,\dot{๐})`$ is $`\sigma `$homogeneous, that is
$$\alpha ,(0,\alpha \dot{๐})=|\alpha |^\sigma (0,\dot{๐})$$
as well as
$$<u>_{\alpha ๐}=\alpha <u>_๐.$$
In particular
$$e_\sigma (\alpha ๐)=e_\sigma (๐),\alpha ,๐\mathrm{\Lambda }_\sigma .$$
(3.10)
In addition, $`e_\sigma (๐)`$ is clearly invariant under diagonal shifts $`๐๐(t)+\beta (t)๐`$ where $`๐=(1,\mathrm{}1)^N`$. Define now
$$\mathrm{\Lambda }_\sigma ^0:=\{๐=(y_1,\mathrm{}y_N)\mathrm{\Lambda }_\sigma ;y_10.\}$$
and
$$S_\sigma =\{๐\mathrm{\Lambda }_\sigma ^0;_0^1|\dot{๐}(t))|^\sigma dt=1\};B_\sigma =\{๐\mathrm{\Lambda }_\sigma ^0;_0^1|(\dot{๐}(t))|^\sigma dt1\}.$$
(3.11)
By the scaling (3.10) and the diagonal shift invariance we conclude that
$$\overline{e}_\sigma :=\underset{๐\mathrm{\Lambda }_\sigma }{sup}e_\sigma (๐)=\underset{๐S_\sigma }{sup}e_\sigma (๐)=\underset{๐B_\sigma \{0\}}{sup}e_\sigma (๐).$$
(3.12)
Let now $`๐_j`$ be a maximizing sequence of $`e_\sigma `$ in $`S_\sigma `$. There is a weak limit $`๐_{\mathrm{}}B_\sigma `$ of this sequence. The inequality
$$\underset{j\mathrm{}}{lim}๐^{๐_j}(0)๐^๐_{\mathrm{}}(0)$$
holds since $`๐^๐(0)`$ is upper-semi-continuous, but we do not have the same claim for $`<u>_๐`$. So, we cannot prove that $`vecy_{\mathrm{}}`$ is a maximizer of $`e_\sigma `$.
Another problem is that we may have $`๐_{\mathrm{}}=\mathrm{๐}`$, so $`e_\sigma (๐_{\mathrm{}})`$ is not defined at all. As an example, let $`๐\mathrm{\Lambda }_\sigma `$ and assume that $`<u>_๐0`$. This, in particular, implies that $`x(t):=^tu\left(\dot{๐}\right)=\stackrel{~}{x}(t)+\lambda t`$ where $`\stackrel{~}{x}`$ is a periodic function and $`\lambda 0`$. If we replace $`๐`$ by $`๐_j=๐_j(t):=j^1๐(jt)`$ for $`j`$, using the periodicity of $`๐`$ to define $`๐_j`$ on $`(0,1)`$, the following claims are straightforward:
a) $`๐_j\mathrm{\Lambda }_\sigma `$ for any $`j`$.
b) $`lim_j\mathrm{}๐_j=0`$ weakly.
c) $`x_j=j^1x(jt)=j^1\stackrel{~}{x}(jt)+\lambda t=^tu\left(\dot{๐_j}\right)`$.
d) $`<u>_{๐_j}=\lambda `$ for any $`j`$.
In particular, we find out that $`e_\sigma (๐_j)=e_\sigma (๐)`$, while $`e_\sigma `$ is not defined for the weak limit $`lim_j\mathrm{}๐_j=\mathrm{๐}`$.
### 3.2 Relaxation
To overcome the last difficulty we shall extend the definition (3.5-3.7) as follows: Let
$$๐ซ^N:=\{ProbabilityBorelmeasureson^N\}$$
Given $`\nu ๐ซ^N`$, let
$$๐^\nu (p):=_^N(p,v)\nu (dv)$$
(3.13)
and
$$<u>_\nu :=_^Nu(v)\nu (dv).$$
(3.14)
The generalization of (3.8) is given by
$$e_\sigma (\nu ):=\frac{\left|<u>_\nu \right|^\sigma }{๐^\nu (0)}=\frac{\left|<u>_\nu \right|^\sigma }{_^n_1^Nm_i|v_i+u(๐)|^\sigma \nu (d๐)}$$
(3.15)
We shall further extend the definitions of $`\mathrm{\Lambda }_\sigma `$ and $`\mathrm{\Lambda }_\sigma ^0`$ as follows:
$$\overline{\mathrm{\Lambda }}_\sigma :=\{\nu =\nu (dv)๐ซ^N;_^N|v|^\sigma \nu (dv)<\mathrm{};_^Nv_i\nu (dv)=0,1iN\}$$
(3.16)
$$\overline{\mathrm{\Lambda }}_\sigma ^0:=\left\{\nu \overline{\mathrm{\Lambda }}_\sigma ;\nu =\delta _{v_1}\mu (dv_2,\mathrm{}dv_N)where\mu ๐ซ^{N1}\right\}$$
(3.17)
###### Lemma 3.3.
For any $`๐ฒ\mathrm{\Lambda }_\sigma `$ (res. $`๐ฒ\mathrm{\Lambda }_\sigma ^0`$) there exists $`\nu \overline{\mathrm{\Lambda }}_\sigma `$ (res. $`\nu \overline{\mathrm{\Lambda }}_\sigma ^0`$) so that $`e_\sigma (\nu )=e_\sigma (๐ฒ)`$. Conversely, for any $`\nu \overline{\mathrm{\Lambda }}_\sigma `$ (res. $`\nu \overline{\mathrm{\Lambda }}_\sigma ^0`$) there exists $`๐ฒ\mathrm{\Lambda }_\sigma `$ (res. $`๐ฒ\mathrm{\Lambda }_\sigma ^0`$) so that $`e_\sigma (\nu )=e_\sigma (๐ฒ)`$.
Remark: The measure $`\nu `$ associated with $`๐`$ is related to Young measure. In general, however, Young measures are used to study the oscillatory behavior of a weak limit of $`๐^{\mathrm{}}`$ sequences (see, e.g., \[E\]).
###### Proof.
For the first part, define $`\nu (dv)=_0^1_1^N\delta _{v_i\dot{y}_i(t)}`$. To elaborate, the measure $`\nu `$ corresponding to $`๐`$ is obtained by its application on test functions $`\varphi C_0\left(^N\right)`$:
$$_^N\varphi (v)\nu (dv)=_0^1\varphi \left(\dot{๐}(t)\right)๐t.$$
If $`๐\mathrm{\Lambda }_\sigma `$ then the above equality also extend to $`\varphi (\dot{๐})=\dot{๐}`$ and $`\varphi (\dot{๐})=|\dot{๐}|^\sigma `$. In particular
$$_^Nv\nu (dv)=_0^1\dot{๐}(t)๐t=0;_^N|v|^\sigma \nu (dv)=_0^1\left|\dot{๐}\right|^\sigma (t)๐t<\mathrm{}.$$
Finally, the equalities
$$๐^๐(0)=๐^\nu (0);<u>_\nu =<u>_๐$$
(3.18)
hold under this identification.
For the second part we use Theorem 2.1 of \[Am\] to observe the following:
For any such $`\nu `$ there exists a Borel function $`T:[0,1]^N`$ which push forward the Lebesgue measure $`dt`$ on $`[0,1]`$ to $`\nu `$. That is, for any test function $`\varphi C_0\left(^N\right)`$,
$$_0^1\varphi (T(t))๐t=_^N\varphi (v)\nu (dv).$$
In fact, Theorem 2.1 of \[Am\] claims the equality between the infimum of Monge and the minimum of the Kantorowich transport plan of probability measures $`\mu _1`$ to $`\mu _2`$ on $`^N`$, provided $`\mu _1`$ has no atoms. This implies, in particular, that the set of Borel mappings transporting $`\mu _1`$ to $`\mu _2`$ is not empty. In our case we use this result under the identification of $`\mu _1`$ with the Lebesgue measure on $`[0,1]`$ (considered as a Hausdorff measure in $`^N`$ for the embedded interval) and $`\mu _2`$ with $`\nu `$.
So, set $`๐(t):=^tT(s)๐s`$. Is is absolutely continues function satisfying $`\dot{๐}=T`$ a.e. The definition (3.16) implies also that $`_0^1\left|\dot{๐}\right|^\sigma =|v|^\sigma \nu (dv)<\mathrm{}`$, as well as $`_0^1\dot{๐}๐t=v\nu (dv)=0`$, which yields the periodicity of $`๐`$. The equality (3.18) holds under this identification as well. โ
### 3.3 Proof of Theorem 2
Define now
$$\psi (๐):=|u(๐)|+\underset{i=2}{\overset{N}{}}|v_i+u(๐)|,๐=(0,v_2,\mathrm{}v_N).$$
(3.19)
We shall summertime some properties of $`\psi `$ which will be needed later:
###### Lemma 3.4.
There exists $`A>0`$ so that
$$A^1|๐||\psi (v)|A|๐|๐=(0,v_2,\mathrm{}v_N).$$
(3.20)
In addition, for any $`\nu \overline{\mathrm{\Lambda }}_\sigma ^0`$,
$$A^1_^N|\psi (v)|^\sigma \nu (dv)<๐^\nu (0)A_^N|\psi (v)|^\sigma \nu (dv).$$
(3.21)
###### Proof.
The estimate (3.20) follows from the homogeneity of $`u`$, namely $`u(\alpha ๐)=\alpha u(๐)`$, as well as from the evident property $`u(0,v,\mathrm{}v)=0v=0`$. The estimate (3.21) follows from (3.3, 3.4) and (3.13), as well as (3.20). โ
Given $`1<q<\sigma `$, let $`\overline{S}_\sigma ^q`$ be the unit sphere (res. $`\overline{B}_\sigma ^q`$ the unit ball) defined by
$$\overline{S}_\sigma ^q=\{\nu \overline{\mathrm{\Lambda }}_\sigma ^0;_{^{N1}}|๐|^q\nu (dv)=1\};\overline{B}_\sigma ^q=\{\nu \overline{\mathrm{\Lambda }}_\sigma ^0;_{^{N1}}|๐|^q\nu (dv)1\}$$
(3.22)
The analogous of (3.12) also holds due to the scaling invariance (3.10):
###### Lemma 3.5.
$`\overline{e}_\sigma =sup_{\nu \overline{S}_\sigma ^q}e_\sigma (\nu ).`$
###### Proof.
We only have to show that any $`\nu \overline{\mathrm{\Lambda }}_\sigma `$ (res. $`\nu \overline{\mathrm{\Lambda }}_\sigma ^0`$) can be transformed into $`\widehat{\nu }\overline{S}_\sigma `$ (res. $`\widehat{\nu }\overline{S}_\sigma ^0`$) so that $`e_\sigma (\nu )=e_\sigma (\widehat{\nu })`$. Let $`\widehat{\nu }(dv)=\beta ^N\nu (\beta dv)`$ for $`\beta >0`$. The homogeneity properties (3.10) imply that, indeed, $`e_\sigma (\widehat{\nu })=e_\sigma (\nu )`$ for any such $`\beta `$. In addition, $`v\widehat{\nu }(dv)=\beta ^1v\nu (dv)=0`$ if $`\nu \overline{\mathrm{\Lambda }}_\sigma `$. However, $`|v|^q\widehat{\nu }(dv)=\beta ^q|v|^q\nu (dv)`$, so $`\widehat{\nu }\overline{S}_\sigma `$ if $`\beta =\left(|v|^q\nu (dv)\right)^{1/q}`$. โ
Let $`\nu _j`$ be a maximizing sequence of $`e_\sigma `$ in $`\overline{S}_\sigma ^0`$. Since the $`q>1`$ moments of $`\nu _j`$ are uniformly bounded, this sequence is compact (tight) in the weak topology of measures. Let $`\nu _{\mathrm{}}`$ be the weak limit of $`\nu _j`$. Since $`|u(v)|A|v|`$ for some $`A>0`$ it follows that the sequence of (signed) measures $`u(v)\nu _j(dv)`$ is tight as well, and that
$$\underset{j\mathrm{}}{lim}u(v)\nu _j(dv)=u(v)\nu _{\mathrm{}}(dv).$$
(3.23)
On the other hand, it is not a-priori evident that $`\nu _{\mathrm{}}\overline{S}_\sigma ^0`$, since the sequence $`|v|^q\nu _j`$ is not necessarily tight, so we only know
$$|v|^q\nu _{\mathrm{}}(dv)1.$$
(3.24)
We claim, however, that, in fact, the $`\sigma `$ moments of $`\nu _j`$ are uniformly bounded. For, if $`|v|^\sigma \nu _j(dv)\mathrm{}`$, then by (3.20, 3.21), also $`๐^{\nu _j}(0)\mathrm{}`$. Since we also now that $`<u>_{\nu _j}`$ are uniformly bounded, then
$$e_\sigma (\nu _j)=\frac{<u>_{\nu _j}}{๐^{\nu _j}(0)}0,$$
contradicting the assumption that $`\nu _j`$ is a maximum sequence for $`e_\sigma `$.
Since $`q<\sigma `$ by assumption it follows that $`|v|^q\nu _j`$ is a tight sequence as well, so there is, in fact, an equality in (3.24). In particular, it follows that $`\nu _{\mathrm{}}\delta _0`$. Moreover, using (3.20, 3.21) for $`\nu _{\mathrm{}}`$ and the equality in (3.24) again, we also have
$$๐^\nu _{\mathrm{}}(0)>0.$$
On the other hand, $`๐^\nu (0)`$ is nothing but the expectation of $`\nu `$ with respect to a positive, continuous function $`(v)=_{i=2}^Nm_i|v_iu(v)|^\sigma `$. Hence
$$\underset{j\mathrm{}}{lim}๐^{\nu _j}(0):=\underset{j\mathrm{}}{lim}(v)\nu _j(dv)(v)\nu _{\mathrm{}}(dv):=๐^\nu _{\mathrm{}}(0).$$
This, together with (3.23), implies that
$$e_\sigma (\nu _{\mathrm{}})\underset{j\mathrm{}}{lim}e_\sigma (\nu _j).$$
(3.25)
Again, $`\nu _j`$ is a maximizing sequence for $`e_\sigma `$ hence there is an equality in (3.25), so $`\overline{e}_\sigma =e_\sigma (\nu _{\mathrm{}})`$. By the second part of Lemma 3.3 we obtain the existence of $`๐_{\mathrm{}}\mathrm{\Lambda }_\sigma ^0`$ for which $`e_\sigma (\nu _{\mathrm{}})=e_\sigma (๐_{\mathrm{}})`$.
References
\[Am\] Ambrosio, L, Lecture notes on Pptimal Transport problems, Lect. Notes in Math., โMathematical aspects of evolving interfacesโ (CIME seris Madeira (PT) 2000, 1812, P. Colli and J.F. Rodrigues, Eds., 1-52, 2003
\[AGK\] J.E. Avron, O ; Gat and O. Kenneth, Optimal swimming at low Reynolds numbers, Phys. Rev. Lett., 93, # 18, 2004.
\[CHK\] Chipot, M ; Hastings, S and Kinderlehrer, D: Transport in a molecular system, M2AN Math. Model. Num. Anal., 38 , #6, 1011-1034, 2004
\[CKK\] Chipot, M; Kinderlehrer, D and Kowalczyk, M: A variational principle for molecular motors, dedicated to Piero Villagio on the occasion of his 60 birthday, Mechanica 38, #5, 505-518, 2003
\[DKK\] Dolbeault, J; Kinderlehrer, D and Kowalczyk, M:Remarks about the flashing rachet, Partial differential equations and inverse problems, Contemp. Math. 362, 167-175, 2004
\[E\] Evans, L.C.: Weak convergence methods for no nonlinear partial differential equations, CBMS, 74, 1990
\[W\] Wolansky, G: Rotation numbers for measure-valued circle maps, J. Dโanalyse, to appear.
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# Ferromagnetism of the Hubbard Model at Strong Coupling in the Hartree-Fock Approximation
## 1 Introduction
The (one-band) Hubbard model has become a standard model for correlated electrons in condensed matter physics since it is, perhaps, the simplest possible model of itinerant interacting electrons. In spite of its simplicity, its zero temperature phase diagram is rich with different magnetic phases such as paramagnetic, ferromagnetic, and antiferromagnetic phases, depending on the details of the hopping amplitudes, the (relative) coupling constant $`U/t`$ and the filling parameter $`\nu =N/(2|\mathrm{\Lambda }|)`$.
As the Hubbard model is a many-body fermion model, the computation of its ground state for large lattices is a difficult, if not impossible, task, except in one-dimension . Thus various schemes have been developed during the past decades to derive an approximate ground state and then to study its magnetic phase diagram.
In the present paper, we consider the Hartree-Fock approximation of the (repulsive, one-band, nearest-neighbor-hopping) Hubbard model with the intention of studying the validity of the Hartree-Fock approximation. We require the Slater determinants entering the Hartree-Fock energy functional to be eigenfunctions of the operator $`๐_z:=_{x\mathrm{\Lambda }}\{n_{x,}n_{x,}\}`$ of total spin in the $`z`$-direction, and for this reason we refer to the model as the HFz approximation. Our requirement means that each orbital has the form $`\phi (x)|`$ or $`\phi (x)|`$. This is a restriction in the sense that general orbitals are of the form $`\phi (x,\sigma )`$, in which the spin direction depends on position. No other restriction is imposed on the variational states; in particular, no assumption about translation invariance is made a priori. For the HFz model, at small chemical potential and for sufficiently strong repulsion, we give a mathematical proof of *saturated ferromagnetism* in the Hartree-Fock ground state. That is, the HF ground state has maximal total spin and maximal ferromagnetic long-range spatial order. The smallness of the chemical potential and the large strength of the repulsion also insure that the HF ground state density is strictly below half-filling.
Before we come to a detailed description of our result and its proof, we discuss it in comparison to other works.
The appearance of ferromagnetic behaviour has been anticipated in many studies of the Hubbard model and approximations thereof. Among these are (restricted) Hartree-Fock approximations , DMFT models in the limit of infinite spatial dimension , exact diagonalizations on small lattices , variational calculations and studies at low filling . These studies support the conjecture that, for large coupling $`U/t1`$ and away from half-filling, $`\nu 1/2`$, the ground state of the Hubbard model is ferromagnetic. Ferromagnetism has been established for the (full) Hubbard model in case the dispersion relation leads to a very high density of states around the Fermi energy and in case of next-nearest-neigbor hopping .
As said before, the main purpose of the present paper is to prove ferromagnetic behaviour with mathematical rigor. None of the papers cited above match the standards of a mathematical proof: The orbitals in the Hartree-Fock approximation are a priori assumed to be composed of only few Fourier modes; the error terms when taking the limit of infinite spatial dimension in DMFT are not under control; exact diagonalizations are restricted to very small lattices and the implication of these to the thermodynamic limit remains unclear. The work by Mielke and Tasaki is mathematically rigrous, but the assumptions made therein about the lattice structure are rather special. On the other hand, by adding next-nearest-neigbor hopping (two-band Hubbard model), Tasaki has found a Hubbard model that displays ferromagnetism in all dimensions. Tasaki also reviews rigorous results on ferromagnetism in the Hubbard model in .
While the prediction of ferromagnetism in the Hubbard model and approximations thereof is supported by the above studies, we also know that HF theory predicts anti-ferromagnetism (in the sense that the total spin is zero) at higher densities, notably at half-filling . Furthermore, our proof shows saturated ferromagnetism at low density and sufficiently large coupling in HF theory, even in one-dimension, but the actual ground state always has spin zero in one-dimension as long as there is only nearest-neighbor hopping (see ).
Even more seriously, our conclusion is opposite to what actually occurs in the Hubbard model. Namely, at very low density (and independent of the value of $`U>0`$), there is no magnetization in the ground state of this model. In the ground state $`๐_z`$ is close to zero and converges to zero, as the particle density tends to zero. This has been pointed out in , based on arguments similar to the following transcription to lattice systems of the recent work .
In this paper it was shown that fermions in the 3-dimensional continuum $`^3`$ (instead of the lattice $`^3`$), and with a repulsive two-body potential, have a ground state energy density, $`e`$, given by
$$e(\rho _{},\rho _{})=\frac{\mathrm{}^2}{2m}\frac{3}{5}(6\pi ^2)^{2/3}\left(\rho _{}^{5/3}+\rho _{}^{5/3}\right)+\frac{\mathrm{}^2}{2m}8\pi a\rho _{}\rho _{}+\mathrm{higher}\mathrm{order}\mathrm{in}(\rho _{},\rho _{}),$$
(1.1)
where $`\rho _{},\rho _{}`$ are the densities of the โspin-upโ and the โspin-downโ fermions and $`a`$ is the scattering length of the two-body potential. Because $`\rho ^{5/3}`$ dominates $`\rho ^2`$ for small $`\rho `$, it is clear from (1.1) that the minimum energy occurs approximately, if not exactly, when $`\rho _{}=\rho _{}=\rho /2`$. This answers the questions in \[21, problem 3\].
To show that there is vanishing net magnetization as $`\rho 0`$ one only needs an upper bound for $`e`$ of the form (1.1). For the Hubbard model (where the two-body potential is a positive delta-function, or even a hard core) this can conveniently be done by a variational wave-function of the form $`\mathrm{\Psi }=F\mathrm{\Psi }_0`$, where $`\mathrm{\Psi }_0`$ is a Slater determinant, and $`F`$ is the projection onto the states with no double occupancy โ in imitation of . We omit the details, but we draw attention to the fact that $`F\mathrm{\Psi }_0`$ is not a Slater determinant, reflecting the more complex structure of correlations in the actual ground state of the Hubbard model. The proof of an analog of (1.1) with precise constants is a more complicated matter which is now under investigation, but it is not needed for the present discussion.
Our setting is the usual (repulsive) Hubbard model with nearest-neighbor hopping on a $`d`$-dimensional cubic lattice $`\mathrm{\Lambda }`$, with periodic boundary conditions and linear size $`L`$, which we assume to be an even integer. It is defined by the second quantized Hamiltonian
$$H_{\mu ,U}=\underset{x,y\mathrm{\Lambda },\sigma =,}{}(\mathrm{\Delta }_{x,y}\mu \delta _{x,y})c_{x,\sigma }^{}c_{y,\sigma }^{}+U\underset{x\mathrm{\Lambda }}{}n_{x,}n_{x,}.$$
(1.2)
We work at fixed chemical potential $`\mu `$ instead of fixed particle number. The only slightly unusual notation is $`\mathrm{\Delta }_{x,y}=T_{x,y}2d\delta _{x,y}`$ for the matrix elements of the discrete Laplacian $`\mathrm{\Delta }`$ on $`\mathrm{\Lambda }`$, with $`T_{x,y}:=1\mathrm{l}[|xy|_1=1]`$ being the nearest-neighbor hopping matrix and $`\delta _{x,y}=1\mathrm{l}[x=y]`$ the Kronecker-Delta.
The operators $`c_{x,\sigma }^{}`$, $`c_{x,\sigma }^{}`$, and $`n_{x,\sigma }:=c_{x,\sigma }^{}c_{x,\sigma }^{}`$ are the usual fermion creation, annihilation, and number operators, respectively, at site $`x\mathrm{\Lambda }`$ and of spin $`\sigma \{,\}`$, obeying the canonical anticommutation relations $`\{c_{x,\sigma }^{},c_{y,\tau }^{}\}=\{c_{x,\sigma }^{},c_{y,\tau }^{}\}`$ $`=0`$, $`\{c_{x,\sigma }^{},c_{y,\tau }^{}\}=\delta _{x,y}\delta _{\sigma ,\tau }`$, and $`c_{x,\sigma }^{}|0=0`$, for all $`x,y,\sigma ,\tau `$. Here $`|0`$ is the vacuum vector in the usual Fock space $`_\mathrm{\Lambda }:=_f(^\mathrm{\Lambda }^2)`$ of spin-$`\frac{1}{2}`$ fermions. The Hamiltonian $`H_{\mu ,U}`$ depends parametrically on the chemical potential $`\mu >0`$ and the coupling constant $`U>0`$.
Note that the usual hopping parameter $`t`$ equals $`1`$ here and that the discrete Laplacian $`\mathrm{\Delta }`$ differs from the usual hopping matrix by the inclusion of the diagonal term, i.e., $`2d`$ times the identity matrix. This difference amounts to a convenient redefinition of the chemical potential $`\mu `$, so that $`\mu =0`$ corresponds precisely to zero filling since the hopping matrix $`\mathrm{\Delta }0`$ is a positive semi-definite matrix. Moreover, the boundedness $`0<\mu <4d`$ of $`\mu `$ together with the assumption that $`U4d`$ insures that the corresponding electron density in the HF ground state is always at low filling, i.e., strictly below half-filling, $`0\rho <1`$.
Our definition of $`\mu `$ is convenient because in this paper, we are concerned with the Hubbard model at low filling, and Our assumption of a bounded chemical potential $`0\mu 2d`$
Apart from this, everything is standard.
The Hamiltonian $`H_{\mu ,U}`$ is a linear operator on the Fock space and the ground state energy $`E_{\mu ,U}^{(\mathrm{gs})}`$ is its smallest eigenvalue,
$$E_{\mu ,U}^{(\mathrm{gs})}:=\mathrm{min}\{\mathrm{\Psi }|H\mathrm{\Psi }|\mathrm{\Psi }_\mathrm{\Lambda },\mathrm{\Psi }=1\}.$$
(1.3)
As the dimension $`dim(_\mathrm{\Lambda })=2^{dim(^\mathrm{\Lambda }^2)}=4^{(L^d)}<\mathrm{}`$ is finite, the determination of $`E_{\mu ,U}^{(\mathrm{gs})}`$ amounts to diagonalizing the finite-dimensional, selfadjoint matrix $`H_{\mu ,U}`$. The fast growth of this dimension with the number $`L^d`$ of points in the lattice $`\mathrm{\Lambda }`$, however, allows for an explicit diagonalization of $`H_{\mu ,U}`$ by a modern computer only up to $`L=4`$, in three spatial dimensions, $`d=3`$.
The Hartree-Fock (HF) approximation is an important method to reduce the high-dimensional many-particle problem given by the diagonalization of $`H_{\mu ,U}`$ to a low-dimensional, but nonlinear variational problem. It is defined by restricting the minimization in (1.3) to Slater determinants $`\phi _1\mathrm{}\phi _N`$, where $`\{\phi _i\}_{i=1}^N^\mathrm{\Lambda }^2`$ is an orthonormal family of $`N`$ one-electron wave functions. The HF approximation to the Hubbard model was analyzed in in the special situation when the number of electrons equals the number of lattice sites, $`N=|\mathrm{\Lambda }|`$, which is usually referred to as *half-filling*.
Note that a priori no other condition but orthonormality is imposed on the orbitals $`\{\phi _i\}_{i=1}^N`$ in the Slater determinants varied over in Hartree-Fock theory. This is sometimes stressed by calling it the *unrestricted Hartree-Fock theory*. Let us temporarily consider a general many-body Hamiltonian $`H`$ which commutes with a certain symmetry operator $`S`$, i.e., $`[H,S]=0`$. It is important to note that in this case, the HF ground state $`\mathrm{\Phi }_{hf}`$, i.e., the Slater determinant which minimizes the energy $`\mathrm{\Phi }_{hf}|H\mathrm{\Phi }_{hf}`$, is not necessarily an eigenstate of $`S`$. Phrased differently, unrestricted Hartree-Fock theory may (depending on the model) break the symmetry $`S`$. The following are examples that occur in physically relevant situations: unrestricted HF ground states of atoms are, in general, not eigenfunctions of the angular momentum operator (because in unrestricted HF theory, all shells are filled ) - even though the atomic Hamiltonian is rotationally invariant; the ground state in the BCS theory of superconductors (which is a variant of HF theory) is not an eigenfunction of the number operator - even though the BCS Hamiltonian preserves the particle number; a HF ground state for the Hubbard model with non-zero spin breaks the invariance of the Hubbard Hamiltonian under global spin rotations; charge density waves (CDW) and spin density waves (SDW) of the Hubbard model are translation invariant only by translation of an *even* number of lattice sites, breaking the (full) translation symmetry the Hubbard Hamiltonian $`H_{\mu ,U}`$ posesses. As it is impossible to predict a priori whether a symmetry of the Hamiltonian is preserved or not, we call all variations of $`\mathrm{\Phi }|H\mathrm{\Phi }`$ over Slater determinants $`\mathrm{\Phi }`$ which fulfill an additional constraint *restricted Hartree-Fock theory*.
In this paper, we consider a restricted Hartree-Fock theory, which we term the *HFz approximation*. The further restriction imposed is that we minimize in (1.3) only over Slater determinants $`\mathrm{\Phi }`$ that are eigenfunctions of the operator $`๐_z:=_{x\mathrm{\Lambda }}\{n_{x,}n_{x,}\}`$ of total spin in the $`z`$-direction. One could rephrase our condition by saying that we do not allow for spiral spin density waves (SSDW; see, e.g., ) in (1.3). Once again, it is customary to employ this restriction in HF calculations without explicitly drawing attention to the fact that this is a restriction. (In mentioned above, however, we dealt with truly unrestricted HF theory.)
More concretely, our HF wave functions have the form
$$\mathrm{\Phi }=\underset{i=1}{\overset{N_{}}{}}c_{}^{}(f_i)\underset{j=1}{\overset{N_{}}{}}c_{}^{}(g_i)|0,$$
(1.4)
where $`c_,^{}(f)=_{x\mathrm{\Lambda }}f(x)c_{x,,}^{}`$, the integers $`N_,`$ are the particle numbers, and where the $`f_i`$ and $`g_i`$ are two families of orthonormal wave functions on the lattice $`\mathrm{\Lambda }`$, i.e., $`f_i|f_j=g_i|g_j=\delta _{i,j}`$, with $`f|g:=_{x\mathrm{\Lambda }}\overline{f(x)}g(x)`$ denoting the usual hermitian scalar product for such functions.
It is convenient to rephrase the HFz approximation in terms of one-particle density matrices, i.e., complex, self-adjoint $`\mathrm{\Lambda }\times \mathrm{\Lambda }`$ matrices whose eigenvalues lie between $`0`$ and $`1`$. To this end, we denote
$$K_\mu :=\mathrm{\Delta }\mu $$
(1.5)
and observe that
$`\mathrm{\Phi }|H\mathrm{\Phi }`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N_{}}{}}}f_i|K_\mu f_i+{\displaystyle \underset{j=1}{\overset{N_{}}{}}}g_j|K_\mu g_j`$ (1.6)
$`+U{\displaystyle \underset{x\mathrm{\Lambda }}{}}\left({\displaystyle \underset{i=1}{\overset{N_{}}{}}}|f_i(x)|^2\right)\left({\displaystyle \underset{j=1}{\overset{N_{}}{}}}|g_j(x)|^2\right).`$
Introducing the one-particle density matrices $`\gamma _,`$ corresponding to $`\mathrm{\Phi }`$ by
$$\gamma _{}:=\underset{i=1}{\overset{N_{}}{}}|f_if_i|\text{and}\gamma _{}:=\underset{i=1}{\overset{N_{}}{}}|g_ig_i|,$$
(1.7)
we observe that $`\gamma _,=\gamma _,^{}=\gamma _,^2`$ are orthogonal projections of dimension $`N_,`$ and that the energy expectation value of the Slater determinant $`\mathrm{\Phi }`$ is given by $`\mathrm{\Phi }|H\mathrm{\Phi }=_{\mu ,U}^{(\mathrm{hfz})}(\gamma _{},\gamma _{})`$, where
$$_{\mu ,U}^{(\mathrm{hfz})}(\gamma _{},\gamma _{}):=\mathrm{Tr}\left\{K_\mu (\gamma _{}+\gamma _{})\right\}+U\underset{x\mathrm{\Lambda }}{}\rho _{}(x)\rho _{}(x),$$
(1.8)
and the diagonal matrix elements $`\rho _,(x):=(\gamma _,)_{x,x}`$ of $`\gamma _,`$ are the one-particle densities of the electron with spin up (โ$``$โ) and spin down (โ$``$โ), respectively.
The symbol โ$`\mathrm{Tr}`$โ denotes the usual trace $`\mathrm{Tr}\{A\}=_{x\mathrm{\Lambda }}A_{x,x}`$ of a complex $`\mathrm{\Lambda }\times \mathrm{\Lambda }`$ matrix $`A=(A_{x,y})_{x,y\mathrm{\Lambda }}`$ with $`A_{x,y}`$. That is, โ$`\mathrm{Tr}`$โ is the trace over the states in $`^\mathrm{\Lambda }`$ of a single spinless particle on the lattice $`\mathrm{\Lambda }`$. It does not include spin states, and it is not the trace over states in Fock space.
Let us note that the particle numbers $`N_,`$ are not determined ab initio. We are in the grand canonical ensemble, so they are determined by the condition that the total energy (1.8) is minimized.
These observations motivate us to define the *HFz energy* by the following variational principle over projections:
$$E_{\mu ,U}^{(\mathrm{hfz})}:=\mathrm{min}\left\{_{\mu ,U}^{(\mathrm{hfz})}(\gamma _{},\gamma _{})|\gamma _,=\gamma _,^{}=\gamma _,^2\right\}.$$
(1.9)
The two sets of orthogonal projections on $`^\mathrm{\Lambda }`$ over which we minimize in (1.9) is not really well-suited for a variational analysis. In particular, they are not convex. An observation in , however, states that, because $`U0`$, we will obtain the same value for the minimum if we vary over the larger set of all one-particle density matrices, $`0\gamma _,1`$, not only over projections. (Recall that a density matrix is a hermitean $`\mathrm{\Lambda }\times \mathrm{\Lambda }`$ matrix $`\gamma `$ whose eigenvalues lie between 0 and 1, i.e., $`0\gamma 1`$, as a matrix inequality.) Our extended $`E_{\mu ,U}^{(\mathrm{hfz})}`$ is then
$$E_{\mu ,U}^{(\mathrm{hfz})}=\mathrm{min}\left\{_{\mu ,U}^{(\mathrm{hfz})}(\gamma _{},\gamma _{})|0\gamma _,1\right\}.$$
(1.10)
The evaluation of $`E_{\mu ,U}^{(\mathrm{hfz})}`$ and the determination of those pairs $`(\gamma _{},\gamma _{})`$ of one-particle density matrices that minimize $`_{\mu ,U}^{(\mathrm{hfz})}`$ is the objective of this paper. Our main result is that, for any $`0<\mu <4d`$, the minimal value of $`_{\mu ,U}^{(\mathrm{hfz})}`$ is attained for the saturated ferromagnet, provided $`U<\mathrm{}`$ is sufficiently large.
###### Theorem 1.1 (Ferromagnetism).
For any $`0<\mu <4d`$, there is a finite length $`L_\mathrm{\#}(\mu )`$ and a finite coupling constant $`U_\mathrm{\#}(\mu )0`$, such that, for all even $`LL_\mathrm{\#}(\mu )`$ and all $`UU_\mathrm{\#}(\mu )`$, the minimal HFz energy is given by the sum of the negative eigenvalues of $`\mathrm{\Delta }\mu `$,
$$E_{\mu ,U}^{(\mathrm{hfz})}=\mathrm{Tr}\left\{[\mathrm{\Delta }\mu ]_{}\right\}.$$
(1.11)
If $`\mu `$ is not an eigenvalue of $`\mathrm{\Delta }`$ and if $`(\gamma _{},\gamma _{})`$ is a minimizer of the HFz functional, i.e., $`0\gamma _,1`$, and $`_{\mu ,U}^{(\mathrm{hfz})}(\gamma _{},\gamma _{})=E_{\mu ,U}^{(\mathrm{hfz})}`$, then
either $`\gamma _{}=\mathrm{\hspace{0.25em}1}\mathrm{l}[\mathrm{\Delta }<\mu ],`$ $`\gamma _{}=\mathrm{\hspace{0.25em}0}`$ (1.12)
or $`\gamma _{}=\mathrm{\hspace{0.25em}0},`$ $`\gamma _{}=\mathrm{\hspace{0.25em}1}\mathrm{l}[\mathrm{\Delta }<\mu ],`$ (1.13)
where $`1\mathrm{l}[\mathrm{\Delta }<\mu ]`$ is the spectral projection of $`\mathrm{\Delta }`$ onto $`(\mathrm{},\mu )`$.
\[With reference to Eq. (1.11) and elsewhere, note that in our notation, $`[X]_{}=\mathrm{min}\{X,\mathrm{\hspace{0.17em}0}\}`$ is negative, whereas elsewhere one often defines $`[X]_{}`$ to be positive, i.e., $`[X]_{}:=\mathrm{max}\{X,\mathrm{\hspace{0.17em}0}\}`$. If $`X`$ is a self adjoint operator then $`[X]_{}`$ denotes the negative part of $`X`$ and $`\mathrm{Tr}[X]_{}`$ is the sum of the negative eigenvalues of $`X`$.\]
Theorem 1.1 is not really as complicated as it looks. It is stated in terms of a length $`L_\mathrm{\#}`$ and coupling constant $`U_\mathrm{\#}`$ in order to make it clear that the state of saturated ferromagnetism is obtained not only asymptotically in the thermodynamic limit and asymptotically as $`U\mathrm{}`$, but it holds for all systems with large, finite interaction and sufficiently large size.
Theorem 1.1 states that, for *any* value of the chemical potential $`\mu (0,4d)`$, the HFz variational principle yields a ferromagnetic minimizer, provided $`U`$ and $`L`$ are chosen sufficiently large (but still finite). A similar statement was proved in \[17, Theorem 4.7\] for $`U=\mathrm{}`$ (which amounts to requiring $`\mathrm{\Phi }|n_{x,}n_{x,}\mathrm{\Phi }=0`$, on every lattice site $`x\mathrm{\Lambda }`$).
At first sight, Theorem 1.1 seems to contradict another fact proved in that the HF minimizer is antiferromagnetic at half-filling. But as the definition of the chemical potential $`\mu `$ in present paper differs from its definition in by $`2d+U`$, the parameter range of the present paper and of never overlap and, hence, there is no contradiction.
As just mentioned, the minimal HF energy and the minimal HFz energy agree in the half-filling case, as shown in . We conjecture that this is also the case for the range of the chemical potential $`\mu (0,4d)`$ and sufficiently large $`U`$, but we do not know how to prove this conjecture. This is a topic for future research.
From Theorem 1.1 we conclude that at small filling there is a phase transition (within the context of HFz theory) from paramagnetism for small $`U`$ to saturated ferromagnetism for large $`U`$. This follows from continuity and the fact that when $`U=0`$ we can find the ground state explicitly and, as is well known, it has $`S=0`$ and is obtained from filling up the Fermi sea for both $``$ and $``$ states.
If $`0<\mu \frac{1}{2}`$ then we can estimate $`L_\mathrm{\#}(\mu )`$ and $`U_\mathrm{\#}(\mu )`$ in Theorem 1.1 more explicitly. For the precise formulation of these estimates, we introduce the following constants,
$`L_{}(\mu )`$ $`:=`$ $`2M_{}(\mu ):=\mathrm{\hspace{0.33em}24}(4d)^2\mu ^2,`$ (1.14)
$`\kappa (\mu )`$ $`:=`$ $`{\displaystyle \frac{\mu ^d}{4^{2d+1}e^dd^d}}\left[1+2\mathrm{ln}(2)(d^1+1)+\mathrm{ln}\left(4d\mu ^1\right)\right]^{2d}`$ (1.15)
$`\alpha _{}(\mu )`$ $`:=`$ $`{\displaystyle \frac{|S^{d1}|\mu ^{(2+d)/2}}{2^{1+d/2}(2\pi )^d(4d)^5}}`$ (1.16)
$`\delta _{}(\mu ,\alpha )`$ $`:=`$ $`\mathrm{min}\{{\displaystyle \frac{\alpha ^2}{(12d)^2}},{\displaystyle \frac{\alpha }{3\mu [4M_{}(\mu )+1]^d}},{\displaystyle \frac{\kappa (\mu )}{2}}\}`$ (1.17)
$`U_{}(\mu ,\alpha )`$ $`:=`$ $`\mathrm{max}\{{\displaystyle \frac{2\mu }{\delta _{}(\mu ,\alpha )}},{\displaystyle \frac{24d^2}{\alpha \delta _{}(\mu ,\alpha )}}\},`$ (1.18)
where $`|S^{d1}|=2\pi ^{d/2}/\mathrm{\Gamma }(d/2)`$ is the measure of the unit sphere in $`^d`$.
###### Theorem 1.2.
For any $`0<\mu \frac{1}{2}`$, Theorem 1.1 holds true with $`L_\mathrm{\#}(\mu ):=L_{}(\mu )`$ and $`U_\mathrm{\#}(\mu ):=U_{}(\mu ,\alpha _{}(\mu ))`$, as defined in (1.14), (1.16), and (1.18).
The explicit form of $`L_{}(\mu )`$, $`\alpha _{}(\mu )`$, and $`U_{}(\mu ,\alpha _{}(\mu ))`$, for a given $`0<\mu \frac{1}{2}`$, in Theorem 1.2 allows us to estimate the actual minimal size of $`L`$ and $`U`$ that guarantees saturated ferromagnetism. The distinction between $`\mu 1/2`$ and $`\mu >1/2`$ is not a fundamental one. It is an artifact of the use in Lemma 3.6 of refs. and , whose methods favored this technical distinction.
Acknowledgements: The authors are grateful to Alessandro Giuliani for very helpful discussions and comments about an earlier version of this paper. They also thank Manfred Salmhofer, Jรผrg Frรถhlich, and Daniel Ueltschi for useful discussions. MT thanks the german student exchange service DAAD for a generous stipend, which supported two thirds of his graduate studies. VB and MT gratefully acknowledge financial support from grant no. HPRN-CT-2002-00277 of the European Union and grant no. Ba 1477/3-3 of the Deutsche Forschungsgemeinschaft. EL gratefully acknowledges support from the Alexander von Humboldt Foundation of a fellowship, the U.S. National Science Foundation, grant no. PHY-0133984, and the hospitality of the Mathematics Departments of the University of Mainz and the Technical University of Berlin. The authors appreciate the careful and helpful work of a referee.
## 2 Proofs of Theorems 1.1 and 1.2
This section contains the proofs of our main results, Theorems 1.1 and 1.2, with the aid of several lemmas which will be proved later in Section 3. Here is a brief outline of the strategy of the proof.
$``$ We first reduce the minimization of $`_{\mu ,U}^{(\mathrm{hfz})}(\gamma _{},\gamma _{})`$ in (1.10) over two one-particle density matrices $`\gamma _{}`$ and $`\gamma _{}`$ to the minimization of an effective energy functional $`\stackrel{~}{}_{\mu ,U}^{(\mathrm{hfz})}(\gamma )`$ which depends only one one-particle density matrix $`\gamma `$. It is given as a sum of two terms, $`\stackrel{~}{}_{\mu ,U}^{(\mathrm{hfz})}(\gamma )=\mathrm{Tr}\{K_\mu \gamma \}+\mathrm{Tr}\{[K_\mu +U\rho ]_{}\}`$, where we recall that $`K_\mu =\mathrm{\Delta }\mu `$.
$``$ Given a trial one-particle density matrix $`\gamma `$ and a small number $`\delta >2\mu U^1`$, we introduce the corresponding particle density $`\rho (x):=\gamma _{x,x}`$ and define the regions $`\mathrm{\Omega }:=\{x|\rho (x)<\delta \}`$ and $`\mathrm{\Omega }^c:=\{x|\rho (x)\delta \}`$ of low and high density onto which we project by $`P_\mathrm{\Omega }=_{x\mathrm{\Omega }}|xx|`$ and $`P_\mathrm{\Omega }^{}=1\mathrm{l}P_\mathrm{\Omega }`$, respectively.
$``$ We then use the fact that $`\gamma `$ is mostly localized in the high density region $`\mathrm{\Omega }^c`$. This leads us to estimate the kinetic energy $`\mathrm{Tr}\{\mathrm{\Delta }P_\mathrm{\Omega }\gamma P_\mathrm{\Omega }\}`$ in $`\mathrm{\Omega }`$ by zero and $`\mathrm{Tr}\{\mathrm{\Delta }P_\mathrm{\Omega }^{}\gamma P_\mathrm{\Omega }^{}\}`$ in $`\mathrm{\Omega }^c`$ by the kinetic energy of the free Fermi gas in $`\mathrm{\Omega }^c`$. The localization error is of order of a small constant times the volume $`|\mathrm{\Omega }|`$ of the boundary of $`\mathrm{\Omega }`$. In Lemma 3.1 we give the exact formulation of the bound which we use to estimate the term $`\mathrm{Tr}\{K_\mu \gamma \}`$ in $`\stackrel{~}{}_{\mu ,U}^{(\mathrm{hfz})}(\gamma )`$.
$``$ For the analysis of the term $`\mathrm{Tr}\{[K_\mu +U\rho ]_{}\}`$ in $`\stackrel{~}{}_{\mu ,U}^{(\mathrm{hfz})}(\gamma )`$, we use the fact that $`\mathrm{\Omega }^c`$ is a classically forbidden region, because $`\mu +U\rho \mu +U\delta \mu `$ in $`\mathrm{\Omega }^c`$. So, as shown in Lemma 3.2, we can replace $`\mathrm{Tr}\{[K_\mu +U\rho ]_{}\}`$ by $`\mathrm{Tr}\{[P_\mathrm{\Omega }(K_\mu +U\rho )P_\mathrm{\Omega }]_{}\}`$, up to localization errors of order of a small constant times $`|\mathrm{\Omega }|`$.
$``$ We then pick a (large, but fixed) number $`M>1`$ and further split up the low density region $`\mathrm{\Omega }`$ into the subset $`\mathrm{\Omega }_1`$ of those points in $`\mathrm{\Omega }`$ that are at most at distance $`2M`$ away from the boundary $`\mathrm{\Omega }`$ and the bulk $`\mathrm{\Omega }_2\mathrm{\Omega }`$ of points of distance $`2M`$ or more to $`\mathrm{\Omega }`$. The contribution of $`\mathrm{\Omega }_1`$ turns out to be negligible because $`\mathrm{\Omega }_1`$ contains at most $`(4M+1)^d|\mathrm{\Omega }|`$ points, and the density is low in $`\mathrm{\Omega }_1\mathrm{\Omega }`$.
$``$ The estimate of the region $`\mathrm{\Omega }_2x`$ then uses the lower bound on the spatial density $`1\mathrm{l}[K_\mu +U\rho <0](x,x)`$ of the projection onto the negative eigenvalues of $`K_\mu +U\rho `$ (actually, $`\stackrel{~}{\rho }`$ instead of $`\rho `$), which we derive in Lemma 3.3
$``$ Adding up the estimates derived so far, we finally observe that $`\stackrel{~}{}_{\mu ,U}^{(\mathrm{hfz})}(\gamma )`$ is bounded below by $`\mathrm{Tr}\{[P_\mathrm{\Omega }K_\mu P_\mathrm{\Omega }]_{}\}+\mathrm{Tr}\{[P_\mathrm{\Omega }^{}K_\mu P_\mathrm{\Omega }^{}]_{}\}\eta |\mathrm{\Omega }|=:Y\eta |\mathrm{\Omega }|`$, where $`\eta >0`$ becomes small when $`U1`$ and $`\delta >0`$ is properly chosen. In Lemma 3.6, we reproduce the result from that $`Y`$ can be estimated from below by $`\mathrm{Tr}\{[K_\mu ]_{}\}+\alpha |\mathrm{\Omega }|`$, where $`\alpha >0`$ depends only on $`\mu `$. In other words, the introduction of a domain wall at $`\mathrm{\Omega }`$ drives up the energy by $`\alpha |\mathrm{\Omega }|`$, which dominates $`\eta |\mathrm{\Omega }|`$, provided $`\eta `$ is small. This establishes that $`\stackrel{~}{}_{\mu ,U}^{(\mathrm{hfz})}(\gamma )\mathrm{Tr}\{[K_\mu ]_{}\}+(\alpha \eta )|\mathrm{\Omega }|`$, which implies the claim.
To carry out the proof in detail, we start with the observation that the minimization over two one-particle density matrices in (1.10) can actually be reduced to the minimization over only one one-particle density matrix. To see this, we observe that
$$\underset{x\mathrm{\Lambda }}{}\rho _{}(x)\rho _{}(x)=\mathrm{Tr}\{\rho _{}\gamma _{}\},$$
(2.1)
where $`\rho _{}`$ acts as a multiplication operator, $`\left(\rho _{}f\right)(x):=\rho _{}(x)f(x)`$. Thus we have
$`E_{\mu ,U}^{(\mathrm{hfz})}`$ $`=`$ $`\underset{0\gamma _{}1}{\mathrm{min}}\left[\mathrm{Tr}\{K_\mu \gamma _{}\}+\underset{0\gamma _{}1}{\mathrm{min}}\left(\mathrm{Tr}\{(K_\mu +U\rho _{})\gamma _{}\}\right)\right]`$ (2.2)
$`=`$ $`\underset{0\gamma _{}1}{\mathrm{min}}\left(\mathrm{Tr}\{K_\mu \gamma _{}\}+\mathrm{Tr}\{[K_\mu +U\rho _{}]_{}\}\right).`$ (2.3)
(Recall that $`K_\mu =\mathrm{\Delta }\mu `$.) In other words, we have done the minimization over $`\gamma _{}`$ in (2.2) by taking $`\gamma _{}`$ to be the projection onto the negative eigenspaces of $`K_\mu +U\rho _{}`$. Thus, as our minimization principle over only one $`\gamma `$, we obtain the following.
$`E_{\mu ,U}^{(\mathrm{hfz})}`$ $`=`$ $`\mathrm{min}\left\{\stackrel{~}{}_{\mu ,U}^{(\mathrm{hfz})}(\gamma )|0\gamma 1\right\},`$ (2.4)
$`\stackrel{~}{}_{\mu ,U}^{(\mathrm{hfz})}(\gamma )`$ $`:=`$ $`\mathrm{Tr}\{K_\mu \gamma \}+\mathrm{Tr}\{[K_\mu +U\rho ]_{}\},`$ (2.5)
where $`\rho (x):=\gamma _{x,x}`$. From now on $`\gamma `$, with $`0\gamma 1`$, is an arbitrary, but fixed one-particle density matrix, for which we bound $`\stackrel{~}{}_{\mu ,U}^{(\mathrm{hfz})}(\gamma )`$ from below. (An upper bound that agrees with Theorem 1.1 is readily obtained simply by choosing the variational function consisting of the unperturbed Fermi sea with all particles spin-up or all spin-down.)
For the next step of the proof we introduce a small number $`\delta >2\mu U^1`$, whose precise value will be chosen in the final step of the proof. Given a one-particle density matrix $`0\gamma 1`$ with corresponding density $`\rho (x):=\gamma _{x,x}`$, we write the lattice $`\mathrm{\Lambda }=\mathrm{\Omega }\mathrm{\Omega }^c`$ as a union of two disjoing subsets of $`\mathrm{\Lambda }`$ in the following way.
$`\mathrm{\Omega }`$ $`:=`$ $`\left\{x\mathrm{\Lambda }\right|\rho (x)<\delta \},`$ (2.6)
$`\mathrm{\Omega }^c`$ $`:=`$ $`\left\{x\mathrm{\Lambda }\right|\rho (x)\delta \}.`$ (2.7)
These are the regions of low and high density, respectively. We define the boundary $`\mathrm{\Omega }`$ of $`\mathrm{\Omega }`$ by
$$\mathrm{\Omega }:=\left\{x\mathrm{\Omega }\right|\mathrm{dist}_1(x,\mathrm{\Omega }^c)=1\},$$
(2.8)
where $`\mathrm{dist}_1(x,A)`$ is the length of (number of bonds in) a shortest path joining $`x`$ and some point in $`yA`$. Another useful notion of distance which we shall use is $`\mathrm{dist}_{\mathrm{}}(x,A)`$, which is defined by the condition that $`2\mathrm{dist}_{\mathrm{}}(x,A)+1`$ is the sidelength of the smallest cube centered at $`x`$ that intersects $`A`$. When $`A`$ is a single point $`y`$ these distances are denoted by $`|xy|_1`$ and $`|xy|_{\mathrm{}}`$.
We define $`P_\mathrm{\Omega }`$, $`P_{\mathrm{\Omega }^c}=P_\mathrm{\Omega }^{}`$, and $`P_\mathrm{\Omega }`$ to be the orthogonal projections onto $`\mathrm{\Omega }`$, $`\mathrm{\Omega }^c`$, and $`\mathrm{\Omega }`$, respectively, where the projection onto an arbitrary set $`A\mathrm{\Lambda }`$ is given by
$$(P_Af)(x)=\{\begin{array}{cc}f(x)& \text{for }xA\text{,}\\ 0& \text{for }xA\text{.}\end{array}$$
(2.9)
We further set
$$\stackrel{~}{\rho }(x):=\{\begin{array}{cc}\rho (x),& \text{for }x\mathrm{\Omega }^c\text{,}\\ \mathrm{min}\{\frac{\mu }{2U},\rho (x)\},& \text{for }x\mathrm{\Omega }\text{,}\end{array}$$
(2.10)
and observe that $`\stackrel{~}{\rho }(x)\rho (x)`$, for all $`x\mathrm{\Lambda }`$, which implies that
$$\stackrel{~}{}_{\mu ,U}^{(\mathrm{hfz})}(\gamma )\mathrm{Tr}\{K_\mu \gamma \}+\mathrm{Tr}\{[K_\mu +U\stackrel{~}{\rho }]_{}\}.$$
(2.11)
For brevity, we define $`M:=M_{}(\mu ):=12(\frac{4d}{\mu })^2`$ and note that, by assumption, $`L`$ obeys $`L2M`$. We further decompose $`\mathrm{\Omega }`$ into two disjoint subsets $`\mathrm{\Omega }_1`$ and $`\mathrm{\Omega }_2`$ defined by
$`\mathrm{\Omega }_1`$ $`:=`$ $`\left\{x\mathrm{\Omega }\right|\mathrm{dist}_{\mathrm{}}(x,\mathrm{\Omega }^c)2M\},`$ (2.12)
$`\mathrm{\Omega }_2`$ $`:=`$ $`\left\{x\mathrm{\Omega }\right|\mathrm{dist}_{\mathrm{}}(x,\mathrm{\Omega }^c)>2M\}.`$ (2.13)
We observe that the $`\mathrm{}^{\mathrm{}}`$-distance of the points in $`\mathrm{\Omega }_1`$ to the boundary $`\mathrm{\Omega }`$ of $`\mathrm{\Omega }`$ is less or equal to $`2M`$, so $`\mathrm{\Omega }_1\mathrm{\Omega }+Q(2M)`$, where $`Q(\mathrm{})=\{\mathrm{},\mathrm{},\mathrm{}\}^d+L^d`$. Hence
$$|\mathrm{\Omega }_1||\mathrm{\Omega }||Q(2M)|=(4M+1)^d|\mathrm{\Omega }|,$$
(2.14)
and therefore
$$\underset{x\mathrm{\Omega }}{}\rho (x)=\underset{x\mathrm{\Omega }_1}{}\rho (x)+\underset{x\mathrm{\Omega }_2}{}\rho (x)(4M+1)^d\delta |\mathrm{\Omega }|+\underset{x\mathrm{\Omega }_2}{}\rho (x),$$
(2.15)
since $`\rho \delta `$ on $`\mathrm{\Omega }`$. Eq. (2.15) and Lemma 3.1 yield
$`\mathrm{Tr}\{K_\mu \gamma \}`$ $``$ $`\mathrm{Tr}\left\{[P_\mathrm{\Omega }^{}K_\mu P_\mathrm{\Omega }^{}]_{}\right\}`$
$`\left(4d\delta ^{1/2}+\mu (4M+1)^d\delta \right)|\mathrm{\Omega }|\mu {\displaystyle \underset{x\mathrm{\Omega }_2}{}}\rho (x).`$
Next, we apply Lemma 3.2 which asserts
$$\mathrm{Tr}\{[K_\mu +U\stackrel{~}{\rho }]_{}\}\mathrm{Tr}\left\{[P_\mathrm{\Omega }(K_\mu +U\stackrel{~}{\rho })P_\mathrm{\Omega }]_{}\right\}\frac{8d^2}{U\delta }|\mathrm{\Omega }|.$$
(2.17)
Denoting by $`\chi :=1\mathrm{l}[P_\mathrm{\Omega }(K_\mu +U\stackrel{~}{\rho })P_\mathrm{\Omega }<0]`$ the orthogonal projection onto the subspace of negative eigenvalues of $`P_\mathrm{\Omega }(K_\mu +U\stackrel{~}{\rho })P_\mathrm{\Omega }`$ and $`\rho _\chi (x):=\chi _{x,x}`$ its diagonal matrix element, we observe that
$`\mathrm{Tr}\left\{[P_\mathrm{\Omega }(K_\mu +U\stackrel{~}{\rho })P_\mathrm{\Omega }]_{}\right\}`$ $`=`$ $`\mathrm{Tr}\left\{P_\mathrm{\Omega }(K_\mu +U\stackrel{~}{\rho })P_\mathrm{\Omega }\chi \right\}`$
$`=`$ $`\mathrm{Tr}\left\{P_\mathrm{\Omega }K_\mu P_\mathrm{\Omega }\chi \right\}+U{\displaystyle \underset{x\mathrm{\Omega }}{}}\rho _\chi (x)\stackrel{~}{\rho }(x).`$
By Lemma 3.3, the density $`\rho _\chi `$ is bounded below on $`\mathrm{\Omega }_2`$ by the universal constant $`\kappa (\mu )>0`$ defined in (3.19). Therefore
$$\mathrm{Tr}\{[K_\mu +U\stackrel{~}{\rho }]_{}\}\mathrm{Tr}\left\{[P_\mathrm{\Omega }K_\mu P_\mathrm{\Omega }]_{}\right\}\frac{8d^2}{U\delta }|\mathrm{\Omega }|+\kappa (\mu )\underset{x\mathrm{\Omega }_2}{}U\stackrel{~}{\rho }(x).$$
(2.19)
Adding up (2) and (2.19), we obtain
$`\stackrel{~}{}_{\mu ,U}^{(\mathrm{hfz})}(\gamma )`$ $``$ $`\mathrm{Tr}\left\{[P_\mathrm{\Omega }K_\mu P_\mathrm{\Omega }]_{}\right\}+\mathrm{Tr}\left\{[P_\mathrm{\Omega }^{}K_\mu P_\mathrm{\Omega }^{}]_{}\right\}`$ (2.20)
$`\left\{4d\delta ^{1/2}+\mu (4M+1)^d\delta +{\displaystyle \frac{8d^2}{U\delta }}\right\}|\mathrm{\Omega }|`$
$`+{\displaystyle \underset{x\mathrm{\Omega }_2}{}}\left\{\kappa (\mu )U\stackrel{~}{\rho }(x)\mu \rho (x)\right\},`$
and Lemma 3.6 further yields
$`\stackrel{~}{}_{\mu ,U}^{(\mathrm{hfz})}(\gamma )\mathrm{Tr}\left\{[K_\mu ]_{}\right\}`$ $``$ $`\left\{\alpha (\mu )4d\delta ^{1/2}\mu (4M+1)^d\delta {\displaystyle \frac{8d^2}{U\delta }}\right\}|\mathrm{\Omega }|`$ (2.21)
$`+{\displaystyle \underset{x\mathrm{\Omega }_2}{}}\left\{\kappa (\mu )U\stackrel{~}{\rho }(x)\mu \rho (x)\right\}.`$
We choose
$$\delta :=\delta _{}(\mu )=\mathrm{min}\{\frac{\alpha (\mu )^2}{(12d)^2},\frac{\alpha (\mu )}{3\mu (4M+1)^d},\frac{\kappa (\mu )}{2}\},$$
(2.22)
and we observe that if
$$UU_{}(\mu ,\alpha (\mu ))=\mathrm{max}\{\frac{2\mu }{\delta _{}(\mu ,\alpha (\mu ))},\frac{24d^2}{\alpha (\mu )\delta _{}(\mu ,\alpha (\mu ))}\}$$
(2.23)
then our choice for $`\delta `$ fulfills the requirement $`\delta >2\mu U^1`$. Moreover, Eqs. (2.22) and (2.23) imply that
$$4d\delta ^{1/2}+\mu (4M+1)^d\delta +\frac{8d^2}{U\delta }\frac{\alpha (\mu )}{3}+\frac{\alpha (\mu )}{3}+\frac{\alpha (\mu )}{3}\alpha (\mu ).$$
(2.24)
We further set $`\mathrm{\Omega }_2^{}:=\{x\mathrm{\Omega }_2|\rho (x)\frac{\mu }{2U}\}`$ and $`\mathrm{\Omega }_2^{\prime \prime }:=\{x\mathrm{\Omega }_2|\frac{\mu }{2U}<\rho (x)\delta \}`$, so $`\mathrm{\Omega }_2`$ is the disjoint union of $`\mathrm{\Omega }_2^{}`$ and $`\mathrm{\Omega }_2^{\prime \prime }`$, and by the definition (2.10) of $`\stackrel{~}{\rho }`$, we have that
$`{\displaystyle \underset{x\mathrm{\Omega }_2}{}}\left\{\kappa (\mu )U\stackrel{~}{\rho }(x)\mu \rho (x)\right\}`$
$``$ $`{\displaystyle \underset{x\mathrm{\Omega }_2^{}}{}}\{\kappa (\mu )U\mu \}\rho (x)+{\displaystyle \underset{x\mathrm{\Omega }_2^{\prime \prime }}{}}{\displaystyle \frac{\mu }{2}}\left\{\kappa (\mu )2\delta \right\}\mathrm{\hspace{0.33em}0},`$
since $`\delta \frac{1}{2}\kappa (\mu )`$ and $`U2\mu /\delta _{}(\mu ,\alpha (\mu ))\mu /\kappa (\mu )`$. Eqs. (2.24) and (2) insure that the right side of (2.21) is nonnegative, which immediately implies Theorem 1.1.
Theorem 1.2 is obtained by substituting the explicit value of $`\alpha (\mu )`$ from (3.60) into (2.23) and using $`L_{}(\mu )`$ from (3.60) . QED
## 3 Auxiliary Lemmas
In this section we state and prove the lemmas used in the proof of Theorems 1.1 and 1.2 in Section 2.
### 3.1 The Region $`\mathrm{\Omega }^c`$ of High Density
In this subsection, we estimate $`\mathrm{Tr}\{K_\mu \gamma \}`$ from below. We are guided by the intuition that $`\gamma `$ is essentially localized on $`\mathrm{\Omega }^c`$.
###### Lemma 3.1.
$$\mathrm{Tr}\{K_\mu \gamma \}\mathrm{Tr}\left\{[P_\mathrm{\Omega }^{}K_\mu P_\mathrm{\Omega }^{}]_{}\right\}\mathrm{\hspace{0.33em}4}d\delta ^{1/2}|\mathrm{\Omega }|\mu \underset{x\mathrm{\Omega }}{}\rho (x).$$
(3.1)
*Proof.* Inserting $`1\mathrm{l}=P_\mathrm{\Omega }+P_\mathrm{\Omega }^{}`$ into $`\mathrm{Tr}\{K_\mu \gamma \}`$, we obtain
$`\mathrm{Tr}\{K_\mu \gamma \}`$ $`=`$ $`\mathrm{Tr}\{K_\mu P_\mathrm{\Omega }^{}\gamma P_\mathrm{\Omega }^{}\}\mathrm{\hspace{0.33em}2}\mathrm{Re}\mathrm{Tr}\{P_\mathrm{\Omega }^{}\mathrm{\Delta }P_\mathrm{\Omega }\gamma \}+\mathrm{Tr}\{K_\mu P_\mathrm{\Omega }\gamma P_\mathrm{\Omega }\}`$
$``$ $`\mathrm{Tr}\left\{[P_\mathrm{\Omega }^{}K_\mu P_\mathrm{\Omega }^{}]_{}\right\}\mathrm{\hspace{0.33em}2}{\displaystyle \underset{x\mathrm{\Omega },y\mathrm{\Omega }^c}{}}\mathrm{\Delta }_{x,y}|\gamma _{y,x}|\mu \mathrm{Tr}\{P_\mathrm{\Omega }\gamma P_\mathrm{\Omega }\}`$
$`=`$ $`\mathrm{Tr}\left\{[P_\mathrm{\Omega }^{}K_\mu P_\mathrm{\Omega }^{}]_{}\right\}\mathrm{\hspace{0.33em}2}{\displaystyle \underset{x\mathrm{\Omega },y\mathrm{\Omega }^c}{}}\mathrm{\Delta }_{x,y}|\gamma _{y,x}|\mu {\displaystyle \underset{x\mathrm{\Omega }}{}}\rho (x),`$
where we use that $`\mathrm{\Delta }0`$, that $`P_\mathrm{\Omega }^{}\mathrm{\Delta }P_\mathrm{\Omega }=P_\mathrm{\Omega }^{}\mathrm{\Delta }P_\mathrm{\Omega }`$, and that $`0\gamma 1`$. The latter also implies that $`\rho (y)=\gamma _{y,y}1`$, for all $`y\mathrm{\Lambda }`$. Thus, if $`x\mathrm{\Omega }`$ and $`y\mathrm{\Omega }^c`$, the Cauchy-Schwarz inequality yields $`|\gamma _{y,x}|\sqrt{\gamma _{y,y}\gamma _{x,x}}\delta ^{1/2}`$. Moreover, if $`x\mathrm{\Omega }`$, $`y\mathrm{\Omega }^c`$, and $`\mathrm{\Delta }_{x,y}0`$, then $`y`$ is a neighbor of $`x`$, and we obtain
$$\underset{x\mathrm{\Omega },y\mathrm{\Omega }^c}{}\mathrm{\Delta }_{x,y}|\gamma _{y,x}|\delta ^{1/2}\underset{x\mathrm{\Omega }}{}\underset{y\mathrm{\Lambda }:|xy|=1}{}=2d\delta ^{1/2}|\mathrm{\Omega }|,$$
(3.3)
which completes the proof of (3.1). QED
### 3.2 Decoupling the High and Low Density Regions
This subsection is devoted to showing that $`\mathrm{Tr}\{[K_\mu +U\stackrel{~}{\rho }]_{}\}`$ essentially agrees with the corresponding eigenvalue sum $`\mathrm{Tr}\{[P_\mathrm{\Omega }(K_\mu +U\stackrel{~}{\rho })P_\mathrm{\Omega }]_{}\}`$ for the operator localized on $`\mathrm{\Omega }`$, the reason being that $`\mathrm{\Omega }^c`$ is a classically forbidden region since $`\mu +U\stackrel{~}{\rho }\frac{1}{2}U\delta >0`$ on $`\mathrm{\Omega }^c`$.
###### Lemma 3.2.
$$\mathrm{Tr}\{[K_\mu +U\stackrel{~}{\rho }]_{}\}\mathrm{Tr}\left\{[P_\mathrm{\Omega }(K_\mu +U\stackrel{~}{\rho })P_\mathrm{\Omega }]_{}\right\}\frac{8d^2}{U\delta }|\mathrm{\Omega }|.$$
(3.4)
*Proof.* We wish to apply of the Feshbach projection method. To this end, we first observe the following quadratic form bound,
$$P_\mathrm{\Omega }^{}(K_{\stackrel{~}{\mu }}+U\stackrel{~}{\rho })P_\mathrm{\Omega }^{}P_\mathrm{\Omega }^{}(U\stackrel{~}{\rho }\stackrel{~}{\mu })P_\mathrm{\Omega }^{}\frac{1}{2}U\delta P_\mathrm{\Omega }^{},$$
(3.5)
for any $`\stackrel{~}{\mu }[0,\mu ]`$, since $`\stackrel{~}{\rho }\delta `$ on $`\mathrm{\Omega }^c`$ and $`\delta 2\mu U^1`$. Thus, $`P_\mathrm{\Omega }^{}(K_{\stackrel{~}{\mu }}+U\stackrel{~}{\rho })P_\mathrm{\Omega }^{}`$ is positive and invertible on $`RanP_\mathrm{\Omega }^{}`$, and moreover, we have that
$$P_\mathrm{\Omega }\mathrm{\Delta }P_\mathrm{\Omega }^{}\left[P_\mathrm{\Omega }^{}(K_{\stackrel{~}{\mu }}+U\stackrel{~}{\rho })P_\mathrm{\Omega }^{}\right]^1P_\mathrm{\Omega }^{}\mathrm{\Delta }P_\mathrm{\Omega }\frac{2}{U\delta }P_\mathrm{\Omega }\mathrm{\Delta }P_\mathrm{\Omega }^{}\mathrm{\Delta }P_\mathrm{\Omega }.$$
(3.6)
For $`y\mathrm{\Omega }^c`$ and $`f^\mathrm{\Lambda }`$, the Cauchy-Schwarz inequality implies that
$`f|P_\mathrm{\Omega }\mathrm{\Delta }\mathrm{\hspace{0.17em}1}\mathrm{l}_y\mathrm{\Delta }P_\mathrm{\Omega }f=|(\mathrm{\Delta }P_\mathrm{\Omega }f)[y]|^2=\left|{\displaystyle \underset{x\mathrm{\Omega },|xy|_1=1}{}}f(x)\right|^2`$
$``$ $`\left({\displaystyle \underset{x\mathrm{\Omega },|xy|_1=1}{}}|f(x)|^2\right)\left({\displaystyle \underset{x\mathrm{\Lambda },|xy|_1=1}{}}1\right)=\mathrm{\hspace{0.33em}\; 2}d{\displaystyle \underset{x\mathrm{\Omega },|xy|_1=1}{}}|f(x)|^2,`$
which, by summing over all $`y\mathrm{\Omega }^c`$, yields
$`f|P_\mathrm{\Omega }\mathrm{\Delta }P_\mathrm{\Omega }^{}\mathrm{\Delta }P_\mathrm{\Omega }f`$ $`=`$ $`{\displaystyle \underset{y\mathrm{\Omega }^c}{}}f|P_\mathrm{\Omega }\mathrm{\Delta }\mathrm{\hspace{0.17em}1}\mathrm{l}_y\mathrm{\Delta }P_\mathrm{\Omega }f`$
$``$ $`2d{\displaystyle \underset{x\mathrm{\Omega }}{}}\left\{|f(x)|^2\left({\displaystyle \underset{y\mathrm{\Lambda },|xy|_1=1}{}}1\right)\right\}`$
$``$ $`4d^2{\displaystyle \underset{x\mathrm{\Omega }}{}}|f(x)|^2=\mathrm{\hspace{0.33em}\; 4}d^2f|P_\mathrm{\Omega }f.`$
(We thank D. Ueltschi for pointing out (3.2)โ(3.2) to us.) We conclude that
$$P_\mathrm{\Omega }\mathrm{\Delta }P_\mathrm{\Omega }^{}\left[P_\mathrm{\Omega }^{}(K_{\stackrel{~}{\mu }}+U\stackrel{~}{\rho })P_\mathrm{\Omega }^{}\right]^1P_\mathrm{\Omega }^{}\mathrm{\Delta }P_\mathrm{\Omega }\frac{8d^2}{U\delta }P_\mathrm{\Omega }.$$
(3.9)
The invertibility of $`P_\mathrm{\Omega }^{}(K_{\stackrel{~}{\mu }}+U\stackrel{~}{\rho }+e)P_\mathrm{\Omega }^{}`$ on $`RanP_\mathrm{\Omega }^{}`$ implies the applicability of the Feshbach map, for any $`e[0,\mu ]`$. I.e., for any $`e[0,\mu ]`$,
$`F(e)`$ $`:=`$ $`F_{P_\mathrm{\Omega }}[K_\mu +e+U\stackrel{~}{\rho }]eP_\mathrm{\Omega }`$
$`=`$ $`P_\mathrm{\Omega }(K_\mu +U\stackrel{~}{\rho })P_\mathrm{\Omega }P_\mathrm{\Omega }\mathrm{\Delta }P_\mathrm{\Omega }^{}\left[P_\mathrm{\Omega }^{}(K_\mu +e+U\stackrel{~}{\rho })P_\mathrm{\Omega }^{}\right]^1P_\mathrm{\Omega }^{}\mathrm{\Delta }P_\mathrm{\Omega }`$
is a well-defined matrix on $`RanP_\mathrm{\Omega }`$, and the isospectrality of the Feshbach map guarantees that $`e[\mu ,0)`$ is a negative eigenvalue of $`K_\mu +U\stackrel{~}{\rho }`$ of multiplicity $`m(e)`$ if and only if $`e`$ is an (nonlinear) eigenvalue of $`F(e)`$, i.e., if the kernel of $`F(e)+e`$, as a subspace of $`RanP_\mathrm{\Omega }`$, has dimension $`m(e)`$. Note that $`F`$ is monotonically increasing, as a quadratic form, in $`e>0`$. In particular,
$$F(e)F(0)P_\mathrm{\Omega }(K_\mu +U\stackrel{~}{\rho })P_\mathrm{\Omega }\frac{8d^2}{U\delta }P_\mathrm{\Omega },$$
(3.11)
additionally taking (3.9) into account.
We claim that, for all $`\lambda (0,\mathrm{})`$, the number of eigenvalues of $`K_\mu +U\stackrel{~}{\rho }`$ below $`\lambda `$ is smaller than the number of negative eigenvalues of $`F(\lambda )+\lambda `$,
$$\mathrm{Tr}\left\{1\mathrm{l}[K_\mu +U\stackrel{~}{\rho }<\lambda ]\right\}\mathrm{Tr}_\mathrm{\Omega }\left\{1\mathrm{l}[F(\lambda )+\lambda <0]\right\},$$
(3.12)
where $`\mathrm{Tr}_\mathrm{\Omega }`$ denotes the trace on $`RanP_\mathrm{\Omega }`$. Both sides of Eq. (3.12) are zero and thus fulfill the claimed inequality, for $`\lambda \mu `$. Assume that (3.12) is violated, for some $`\lambda (0,\mathrm{})`$, i.e., that $`\lambda _{}:=inf\{\lambda (0,\mathrm{})|\text{Eq. (}\text{3.12}\text{) holds true}\}>0`$. We show that this assumption leads to a contradiction. Obviously, $`\lambda _{}`$ must be an eigenvalue of $`K_\mu +U\stackrel{~}{\rho }`$, and hence also of $`F(\lambda _{})`$, of multiplicity $`m(\lambda _{})1`$, because only then the left or the right side of (3.12) changes (increases, in fact). Moreover, Eq. (3.12) holds true for $`\lambda =\lambda _{}`$ itself, i.e., the infimum in the definition of $`\lambda _{}`$ is a minimum. Hence, for all sufficiently small $`\epsilon >0`$, the definition of $`\lambda _{}`$ and the monotony of $`F(e)`$ in $`e`$ yield
$`\mathrm{Tr}\left\{1\mathrm{l}[K_\mu +U\stackrel{~}{\rho }<\lambda _{}]\right\}`$ $``$ $`\mathrm{Tr}_\mathrm{\Omega }\left\{1\mathrm{l}[F(\lambda _{})+\lambda _{}<0]\right\}`$ (3.13)
$`\mathrm{Tr}\left\{1\mathrm{l}[K_\mu +U\stackrel{~}{\rho }<\lambda _{}+\epsilon ]\right\}`$ $`>`$ $`\mathrm{Tr}_\mathrm{\Omega }\left\{1\mathrm{l}[F(\lambda _{}\epsilon )+\lambda _{}\epsilon <0]\right\}`$ (3.14)
$``$ $`\mathrm{Tr}_\mathrm{\Omega }\left\{1\mathrm{l}[F(\lambda _{})+\lambda _{}\epsilon <0]\right\}.`$
Choosing $`\epsilon >0`$ so small that $`\lambda _{}`$ is the only eigenvalue of $`K_\mu +U\stackrel{~}{\rho }`$ in the interval $`[\lambda _{},\lambda _{}+\epsilon ]`$, we hence obtain
$`m(\lambda _{})`$ $`=`$ $`\mathrm{Tr}\left\{1\mathrm{l}[0K_\mu +U\stackrel{~}{\rho }+\lambda _{}<\epsilon ]\right\}`$ (3.15)
$`=`$ $`\mathrm{Tr}\left\{1\mathrm{l}[K_\mu +U\stackrel{~}{\rho }<\lambda _{}+\epsilon ]\right\}\mathrm{Tr}\left\{1\mathrm{l}[K_\mu +U\stackrel{~}{\rho }<\lambda _{}]\right\}`$
$`>`$ $`\mathrm{Tr}_\mathrm{\Omega }\left\{1\mathrm{l}[F(\lambda _{})+\lambda _{}<\epsilon ]\right\}\mathrm{Tr}_\mathrm{\Omega }\left\{1\mathrm{l}[F(\lambda _{})+\lambda _{}<0]\right\}`$
$`=`$ $`\mathrm{Tr}_\mathrm{\Omega }\left\{1\mathrm{l}[0F(\lambda _{})+\lambda _{}<\epsilon ]\right\}=m(\lambda _{}),`$
arriving at a contradiction, which proves (3.12), for all $`\lambda (0,\mathrm{})`$. From (3.12) and (3.11), we finally conclude
$`\mathrm{Tr}\{[K_\mu +U\stackrel{~}{\rho }]_{}\}`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}\mathrm{Tr}\left\{1\mathrm{l}[K_\mu +U\stackrel{~}{\rho }<\lambda ]\right\}๐\lambda `$
$``$ $`{\displaystyle _0^{\mathrm{}}}\mathrm{Tr}_\mathrm{\Omega }\left\{1\mathrm{l}[F(\lambda )+\lambda <0]\right\}.`$
$``$ $`{\displaystyle _0^{\mathrm{}}}\mathrm{Tr}_\mathrm{\Omega }\left\{1\mathrm{l}[F(0)+\lambda <0]\right\}`$
$`=`$ $`\mathrm{Tr}\{[F(0)]_{}\}=\mathrm{Tr}_\mathrm{\Omega }\{[F(0)]_{}\}`$
$``$ $`\mathrm{Tr}\left\{[P_\mathrm{\Omega }(K_\mu +U\stackrel{~}{\rho })P_\mathrm{\Omega }]_{}\right\}{\displaystyle \frac{8d^2}{U\delta }}\mathrm{Tr}\{P_\mathrm{\Omega }\}.`$
which is the assertion of Lemma 3.2. QED
### 3.3 The Electron Density in the Bulk
In this subsection we consider the spectral projection
$$\chi :=1\mathrm{l}\left[P_\mathrm{\Omega }(K_\mu +U\stackrel{~}{\rho })P_\mathrm{\Omega }<0\right]=1\mathrm{l}\left[P_\mathrm{\Omega }(\mathrm{\Delta }\mu +U\stackrel{~}{\rho })P_\mathrm{\Omega }<0\right]$$
(3.17)
of $`P_\mathrm{\Omega }(\mathrm{\Delta }\mu +U\stackrel{~}{\rho })P_\mathrm{\Omega }`$ onto its negative eigenvalues. Writing $`\mathrm{\Delta }_\mathrm{\Omega }:=P_\mathrm{\Omega }\mathrm{\Delta }P_\mathrm{\Omega }`$, i.e., $`(\mathrm{\Delta }_\mathrm{\Omega })_{x,y}=\mathrm{\Delta }_{x,y}`$, for $`x,y\mathrm{\Omega }`$, and $`=0`$, otherwise, and $`V_{x\mathrm{\Omega }}V(x)1\mathrm{l}_x:=\mu P_\mathrm{\Omega }U\stackrel{~}{\rho }P_\mathrm{\Omega }`$, we have that
$$\chi =1\mathrm{l}[\mathrm{\Delta }_\mathrm{\Omega }V<0]\text{and}x\mathrm{\Omega }:\frac{1}{2}\mu V(x)\mu ,$$
(3.18)
due to the definition (2.9) of $`\stackrel{~}{\rho }`$. Naive semiclassical intuition tells us that, for $`x\mathrm{\Omega }`$, the particle density $`\rho _\chi (x):=\chi _{x,x}`$ corresponding to the one-particle density matrix $`\chi `$ should be bounded below by the particle density of the Fermi gas given by the one-particle density matrix $`1\mathrm{l}[\mathrm{\Delta }<\mu /2]`$. The purpose of this subsection is to prove such a bound (up to a constant factor) where it can be expected to hold, namely, for those points $`x`$ that are sufficiently far away from the boundary of $`\mathrm{\Omega }`$.
###### Lemma 3.3.
Let $`0<\mu 4d`$, define $`M:=M_{}:=12(\frac{4d}{\mu })^2`$. Suppose that $`L`$ obeys $`L2M`$ and that $`x\mathrm{\Omega }`$, with $`\mathrm{dist}_{\mathrm{}}(x,\mathrm{\Omega })>2M`$. Then
$$\rho _\chi (x)\kappa (\mu ):=\frac{\mu ^d}{4^{2d+1}e^dd^d}\left[1+2\mathrm{ln}(2)\left(d^1+1\right)+\mathrm{ln}\left(4d\mu ^1\right)\right]^{2d}.$$
(3.19)
*Proof.* For any $`\beta >0`$, we note that the map $`^\mathrm{\Omega }`$, $`W(e^{\beta (\mathrm{\Delta }_\mathrm{\Omega }W)})_{x,x}`$ is monotonically increasing in $`W`$. Namely, as $`T_\mathrm{\Omega }=P_\mathrm{\Omega }TP_\mathrm{\Omega }`$ has nonnegative matrix elements, so does $`e^{\epsilon \mathrm{\Delta }_\mathrm{\Omega }}`$,
$$\left(e^{\epsilon \mathrm{\Delta }_\mathrm{\Omega }}\right)_{w,z}=e^{2d\epsilon }\left(e^{\epsilon T_\mathrm{\Omega }}\right)_{w,z}=e^{2d\epsilon }\underset{k=0}{\overset{\mathrm{}}{}}\frac{\epsilon ^k}{k!}\left(T_\mathrm{\Omega }^k\right)_{w,z}0,$$
(3.20)
for all $`w,z\mathrm{\Omega }`$. So, if $`n`$ is an integer and $`W,\stackrel{~}{W}^\mathrm{\Omega }`$ with $`W(z)\stackrel{~}{W}(z)`$, for all $`z\mathrm{\Omega }`$, then we have that
$`\left(\left[e^{\beta \mathrm{\Delta }_\mathrm{\Omega }/n}e^{\beta W/n}\right]^n\right)_{z_0,z_n}`$ $`=`$ $`{\displaystyle \underset{z_1,\mathrm{},z_{n1}\mathrm{\Omega }}{}}\left\{{\displaystyle \underset{j=1}{\overset{n}{}}}\left(e^{\beta \mathrm{\Delta }_\mathrm{\Omega }/n}\right)_{z_{j1},z_j}e^{\beta W(z_j)/n}\right\}`$ (3.21)
$``$ $`{\displaystyle \underset{z_1,\mathrm{},z_{n1}\mathrm{\Omega }}{}}\left\{{\displaystyle \underset{j=1}{\overset{n}{}}}\left(e^{\beta \mathrm{\Delta }_\mathrm{\Omega }/n}\right)_{z_{j1},z_j}e^{\beta \stackrel{~}{W}(z_j)/n}\right\}`$
$`=`$ $`\left(\left[e^{\beta \mathrm{\Delta }_\mathrm{\Omega }/n}e^{\beta \stackrel{~}{W}/n}\right]^n\right)_{z_0,z_n},`$
for all $`z_0,z_n\mathrm{\Omega }`$. Setting $`z_0:=z_n:=x\mathrm{\Omega }`$ and taking the limit $`n\mathrm{}`$, the Lie-Trotter product formula and Eq. (3.21) imply that
$$\left(e^{\beta (\mathrm{\Delta }_\mathrm{\Omega }W)}\right)_{x,x}\left(e^{\beta (\mathrm{\Delta }_\mathrm{\Omega }\stackrel{~}{W})}\right)_{x,x},$$
(3.22)
indeed. In particular,
$$e^{\beta \mu /2}\left(e^{\beta \mathrm{\Delta }_\mathrm{\Omega }}\right)_{x,x}\left(e^{\beta (\mathrm{\Delta }_\mathrm{\Omega }V)}\right)_{x,x},$$
(3.23)
since $`V\frac{1}{2}\mu `$ on $`\mathrm{\Omega }`$. On the other hand, $`\mathrm{\Delta }_\mathrm{\Omega }V\mu `$ and $`\chi ^{}(\mathrm{\Delta }_\mathrm{\Omega }V)\chi ^{}0`$, as quadratic forms. The spectral theorem thus implies that
$`\chi e^{\beta (\mathrm{\Delta }_\mathrm{\Omega }V)}\chi `$ $``$ $`\chi e^{\beta \mu }\chi =e^{\beta \mu }\chi ,`$ (3.24)
$`\chi ^{}e^{\beta (\mathrm{\Delta }_\mathrm{\Omega }V)}\chi ^{}`$ $``$ $`\chi ^{}P_\mathrm{\Omega }.`$ (3.25)
Putting together (3.23), (3.24), and (3.25), using that $`\chi `$ and $`\mathrm{\Delta }_\mathrm{\Omega }V`$ commute, we arrive at
$`e^{\beta \mu /2}\left(e^{\beta \mathrm{\Delta }_\mathrm{\Omega }}\right)_{x,x}`$ $``$ $`\left(e^{\beta (\mathrm{\Delta }_\mathrm{\Omega }V)}\right)_{x,x}`$ (3.26)
$`=`$ $`\left(\chi e^{\beta (\mathrm{\Delta }_\mathrm{\Omega }V)}\chi \right)_{x,x}+\left(\chi ^{}e^{\beta (\mathrm{\Delta }_\mathrm{\Omega }V)}\chi ^{}\right)_{x,x}`$
$``$ $`e^{\beta \mu }\chi _{x,x}+\mathrm{\hspace{0.33em}1}.`$
Solving for $`\rho _\chi (x)=\chi _{x,x}`$, we therefore have
$$\rho _\chi (x)e^{\beta \mu /2}\left[(e^{\beta \mathrm{\Delta }_\mathrm{\Omega }})_{x,x}e^{\beta \mu /2}\right],$$
(3.27)
for any $`x\mathrm{\Omega }`$ and any $`\beta >0`$.
Next, recall that $`Q(M)=\{M,\mathrm{},M\}^d+L^d=\{y\mathrm{\Lambda }:|y|_{\mathrm{}}M\}`$ is the box of sidelength $`2M+1`$ centered at $`0\mathrm{\Lambda }`$. Since $`\mathrm{dist}_{\mathrm{}}(x,\mathrm{\Omega })>2M`$, by assumption, we have that
$$Q(M)z+x\mathrm{\Omega },$$
(3.28)
for all $`zQ(M)`$. By Lemma 3.4, this inclusion implies that
$$\left(\mathrm{exp}[\beta \mathrm{\Delta }_\mathrm{\Omega }]\right)_{x,x}\left(\mathrm{exp}[\beta \mathrm{\Delta }_{Q(M)z+x}]\right)_{x,x}=\left(\mathrm{exp}[\beta \mathrm{\Delta }_{Q(M)}]\right)_{z,z},$$
(3.29)
and by averaging this inequality over $`zQ(M)`$, we obtain
$$\left(\mathrm{exp}[\beta \mathrm{\Delta }_\mathrm{\Omega }]\right)_{x,x}\frac{1}{|Q(M)|}\underset{zQ(M)}{}\left(\mathrm{exp}[\beta \mathrm{\Delta }_{Q(M)}]\right)_{z,z}.$$
(3.30)
Now, we apply Lemma 3.5 and arrive at
$`{\displaystyle \frac{1}{|Q(M)|}}{\displaystyle \underset{zQ(M)}{}}\left(\mathrm{exp}[\beta \mathrm{\Delta }_{Q(M)}]\right)_{z,z}{\displaystyle \frac{e^{d\beta /M}}{(2\pi )^d}}{\displaystyle _{[\pi ,\pi ]^d}}\mathrm{exp}[\beta \omega (k)]d^dk`$ (3.31)
$`=`$ $`\left({\displaystyle \frac{2e^{\beta /M}}{\pi }}{\displaystyle _0^{\pi /2}}\mathrm{exp}[4\beta \mathrm{sin}^2(t)]๐t\right)^d,`$
where $`\omega (k)=\omega (k)=_{\nu =1}^d2\left\{1\mathrm{cos}(k_\nu )\right\}=_{\nu =1}^d4\mathrm{sin}^2(k_\nu /2)`$. Choosing $`\beta 1`$, we observe that $`\frac{1}{\pi }_0^{\sqrt{\beta }\pi }e^{t^2}๐t\frac{1}{\pi }_0^\pi e^{t^2}๐t=\frac{1}{2\sqrt{\pi }}\mathrm{erf}[\pi ]\frac{1}{4}`$. Using this and $`\mathrm{sin}^2(t)t^2`$, we have the following estimate,
$$\frac{2e^{\beta /M}}{\pi }_0^{\pi /2}\mathrm{exp}[4\beta \mathrm{sin}^2(t)]๐t\frac{e^{\beta /M}}{\beta ^{1/2}}\frac{1}{\pi }_0^{\sqrt{\beta }\pi }e^{t^2}๐t\frac{e^{\beta /M}}{4\beta ^{1/2}}.$$
(3.32)
Inserting this estimate into (3.31) and then the result in (3.30) and (3.27), we obtain, for any $`\beta 1`$, that
$$\rho _\chi (x)e^{\beta \mu /2}\left[\frac{e^{d\beta /M}}{4^d\beta ^{d/2}}e^{\beta \mu /2}\right]=e^{\tau d}\left[\left(e^{12d\tau /(M\mu )}\right)^{d/2}\left(\frac{\mu }{16ed}\frac{e^\tau }{\tau }\right)^{d/2}\mathrm{\hspace{0.25em}1}\right],$$
(3.33)
where $`\tau :=\beta \mu /d`$. Note that if we require $`\tau 4`$ then $`\beta =\tau d/\mu 1`$, since $`\mu 4d`$. We may thus replace $`\beta [1,\mathrm{})`$ by $`\tau [4,\mathrm{})`$. Our goal is to choose $`\tau `$ such that
$`\left({\displaystyle \frac{\mu }{16ed}}{\displaystyle \frac{e^\tau }{\tau }}\right)^{d/2}2`$ (3.34)
$`\tau \mathrm{ln}(\tau )Y:=\mathrm{\hspace{0.33em}1}+2\mathrm{ln}(2)\left({\displaystyle \frac{1}{d}}+1\right)+\mathrm{ln}\left({\displaystyle \frac{4d}{\mu }}\right).`$ (3.35)
Note that, due to $`\mu 4d`$,
$$2.381+2\mathrm{ln}(2)Y3\mathrm{ln}\left(16d\mu ^1\right).$$
(3.36)
We choose $`\tau :=Y+2\mathrm{ln}(Y)`$ and observe that $`Y2.38`$ insures $`\tau 4.114`$, as required. Moreover, with this choice, we have
$`\tau \mathrm{ln}(\tau )Y`$ $`=`$ $`2\mathrm{ln}(Y)\mathrm{ln}\left[Y+2\mathrm{ln}(Y)\right]`$
$``$ $`\mathrm{ln}(Y)\mathrm{ln}\left[1+2\mathrm{ln}(Y)Y^1\right]`$
$``$ $`\mathrm{ln}(Y)\mathrm{\hspace{0.25em}2}\mathrm{ln}(Y)Y^1=2\mathrm{ln}(Y)\left({\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{Y}}\right)>0,`$
using that $`\mathrm{ln}(1+\epsilon )\epsilon `$, for $`\epsilon 0`$, and $`Y2.38>2`$. Thus, (3.35) and (3.34) hold true. Additionally, we observe that $`Y3\mathrm{ln}(\frac{16d}{\mu })`$ and
$$\tau Y\underset{r>0}{\mathrm{max}}\left\{1+2\left(\frac{\mathrm{ln}r}{r}\right)\right\}=(1+2/e)Y2Y$$
(3.38)
insures that $`\frac{2d\tau }{\mu }\frac{12d}{\mu }\mathrm{ln}(\frac{16d}{\mu })12(\frac{4d}{\mu })^2M_{}M`$. This, in turn, yields
$$\mathrm{exp}\left[1\frac{2d\tau }{M\mu }\right]1,$$
(3.39)
and by inserting (3.39) and (3.34) into (3.33), we arrive at
$$\rho _\chi (x)e^{\tau d}=\frac{\mu ^d}{4^{2d+1}e^dd^d}\left[1+2\mathrm{ln}(2)\left(d^1+1\right)+\mathrm{ln}\left(4d\mu ^1\right)\right]^{2d}.$$
(3.40)
QED
###### Lemma 3.4.
Let $`A,B\mathrm{\Lambda }`$, with $`AB`$, and denote $`\mathrm{\Delta }_A:=P_A\mathrm{\Delta }P_A`$ and $`\mathrm{\Delta }_B:=P_B\mathrm{\Delta }P_B`$. For all $`xA`$ and all $`\beta >0`$,
$$\left(\mathrm{exp}[\beta \mathrm{\Delta }_A]\right)_{x,x}\left(\mathrm{exp}[\beta \mathrm{\Delta }_B]\right)_{x,x}.$$
(3.41)
*Proof.* We first define the nearest-neighbor hopping matrix $`T`$ on $`\mathrm{\Lambda }`$ by $`T_{w,z}:=1`$ if $`|wz|_1=1`$ and $`T_{w,z}:=0`$, otherwise. For a given subset $`C\mathrm{\Lambda }`$, the matrix $`T_C:=P_CTP_C`$ denotes the hopping matrix restricted to $`C`$. Note that $`\mathrm{\Delta }_C=T_C2dP_C`$ is the difference of the two commuting matrices $`T_C`$ and $`2dP_C`$. Hence, for $`xC`$,
$$\left(\mathrm{exp}[\beta \mathrm{\Delta }_C]\right)_{x,x}=\left(\mathrm{exp}[\beta T_C]\mathrm{exp}[2d\beta P_C]\right)_{x,x}=e^{2d\beta }\left(\mathrm{exp}[\beta T_C]\right)_{x,x}.$$
(3.42)
Due to this identity and the fact that $`xAB`$, Eq. (3.41) is equivalent to
$$\left(\mathrm{exp}[\beta T_A]\right)_{x,x}\left(\mathrm{exp}[\beta T_B]\right)_{x,x}.$$
(3.43)
Now, $`0(T_A)_{w,z}(T_B)_{w,z}`$, and hence $`(T_A^n)_{x,x}(T_B^n)_{x,x}`$, for all intergers $`n`$. Thus, (3.43) follows from an expansion of the exponentials in Taylor series,
$$\left(\mathrm{exp}[\beta T_A]\right)_{x,x}=\underset{n=0}{\overset{\mathrm{}}{}}\frac{\beta ^n}{n!}(T_A^n)_{x,x}\underset{n=0}{\overset{\mathrm{}}{}}\frac{\beta ^n}{n!}(T_B^n)_{x,x}=\left(\mathrm{exp}[\beta T_B]\right)_{x,x}.$$
(3.44)
QED
###### Lemma 3.5.
Let $`Q=\{m,\mathrm{},m\}^d^d`$ be a cube. Denote by $`\mathrm{\Delta }_Q`$ the nearest-neighbor Laplacian on $`Q`$, i.e., $`\mathrm{\Delta }_Q=P_Q\mathrm{\Delta }P_Q=2dP_Q+T_Q`$, $`T_Q:=P_QTP_Q`$, and $`T_{x,y}=1\mathrm{l}(|xy|_1=1)`$. Then, for all $`\beta >0`$,
$$\frac{1}{|Q|}\underset{zQ}{}\left(\mathrm{exp}[\beta \mathrm{\Delta }_Q]\right)_{z,z}\frac{e^{d\beta /m}}{(2\pi )^d}_{[\pi ,\pi ]^d}\mathrm{exp}[\beta \omega (k)]d^dk,$$
(3.45)
where $`\omega (k):=_{\nu =1}^d2\left\{1\mathrm{cos}(k_\nu )\right\}`$.
*Proof.* We may pick an even integer $`r`$, choose $`L:=r(2m+1)`$, and identify $`Q`$ with $`Q+L^d\mathrm{\Lambda }`$. (Note that the statement of the lemma makes no reference to the Hubbard model analyzed before, and for the purpose of the proof, $`L`$ can be taken an arbitrarily large integer multiple of $`2m+1`$.) Given $`s_r^d`$, we define $`Q(s):=Q+(2m+1)s`$ and observe that the family $`\{Q(s)\}_{s_r^d}`$ of cubes define a disjoint partition of $`\mathrm{\Lambda }`$, i.e.,
$$\mathrm{\Lambda }=\underset{s_r^d}{}Q(s)\text{and}ss^{}:Q(s)Q(s^{})=\mathrm{}.$$
(3.46)
Hence
$$\widehat{\mathrm{\Delta }}:=\underset{s_r^d}{}\mathrm{\Delta }_{Q(s)}$$
(3.47)
is the sum of translated, but mutually disconnected copies of $`\mathrm{\Delta }_Q`$. We observe that
$`\mathrm{Tr}\left\{\mathrm{exp}[\beta \widehat{\mathrm{\Delta }}]\right\}`$
$`=`$ $`{\displaystyle \underset{x\mathrm{\Lambda }}{}}\left(\mathrm{exp}[\beta \widehat{\mathrm{\Delta }}]\right)_{x,x}={\displaystyle \underset{s_r^d}{}}{\displaystyle \underset{zQ}{}}\left(\mathrm{exp}[\beta \widehat{\mathrm{\Delta }}]\right)_{z+(2m+1)s,z+(2m+1)s}`$
$`=`$ $`{\displaystyle \underset{s_r^d}{}}{\displaystyle \underset{zQ}{}}\left(\mathrm{exp}[\beta \mathrm{\Delta }_{Q(s)}]\right)_{z+(2m+1)s,z+(2m+1)s}=r^d{\displaystyle \underset{zQ}{}}\left(\mathrm{exp}[\beta \mathrm{\Delta }_Q]\right)_{z,z}.`$
As an intermediate result, we thus have
$$\frac{1}{|Q|}\underset{zQ}{}\left(\mathrm{exp}[\beta \mathrm{\Delta }_Q]\right)_{z,z}=\frac{1}{|\mathrm{\Lambda }|}\mathrm{Tr}\left\{\mathrm{exp}[\beta \widehat{\mathrm{\Delta }}]\right\},$$
(3.49)
since $`|\mathrm{\Lambda }|=L^d=r^d|Q|`$.
Next, we translate $`\widehat{\mathrm{\Delta }}`$ by the elements of $`Q`$, i.e., for $`\eta Q`$, we introduce $`\widehat{\mathrm{\Delta }}^{(\eta )}`$ on $`^\mathrm{\Lambda }`$ by
$$\widehat{\mathrm{\Delta }}^{(\eta )}:=\underset{q_r^d}{}\mathrm{\Delta }_{Q(q)+\eta }=\underset{q_r^d}{}\mathrm{\Delta }_{Q+\eta +(2m+1)q}.$$
(3.50)
Of course, $`\widehat{\mathrm{\Delta }}^{(\eta )}`$ is unitarily equivalent to $`\widehat{\mathrm{\Delta }}`$. We observe that
$$\frac{1}{|Q|}\underset{\eta Q}{}\widehat{\mathrm{\Delta }}^{(\eta )}=\frac{1}{|Q|}\underset{y\mathrm{\Lambda }}{}\mathrm{\Delta }_{Q+y}=2d1\mathrm{l}_^\mathrm{\Lambda }+\frac{1}{|Q|}\underset{y\mathrm{\Lambda }}{}T_{Q+y},$$
(3.51)
where, for $`w,z\mathrm{\Lambda }`$,
$`\left({\displaystyle \underset{y\mathrm{\Lambda }}{}}T_{Q+y}\right)_{w,z}={\displaystyle \underset{y\mathrm{\Lambda }}{}}1\mathrm{l}_Q(wy)\mathrm{\hspace{0.17em}1}\mathrm{l}_Q(zy)T_{w,z}`$
$`=`$ $`\left|(Q+w)(Q+z)\right|T_{w,z}=\mathrm{\hspace{0.33em}2}m(2m+1)^{d1}T_{w,z},`$
since $`T_{w,z}0`$ only if $`wz`$ are neighboring lattice sites. Hence,
$`{\displaystyle \frac{1}{|Q|}}{\displaystyle \underset{\eta Q}{}}\widehat{\mathrm{\Delta }}^{(\eta )}`$ $`=`$ $`2d1\mathrm{l}_^\mathrm{\Lambda }+{\displaystyle \frac{2m}{2m+1}}T={\displaystyle \frac{2d}{2m+1}}1\mathrm{l}_^\mathrm{\Lambda }+{\displaystyle \frac{2m}{2m+1}}\mathrm{\Delta }`$ (3.53)
$``$ $`{\displaystyle \frac{d}{m}}1\mathrm{l}_^\mathrm{\Lambda }+\mathrm{\Delta }`$
where $`\mathrm{\Delta }0`$ is the nearest-neighbor Laplacian on $`\mathrm{\Lambda }`$ (with periodic b.c.). This and the convexity of $`A\mathrm{Tr}\{e^{\beta A}\}`$ therefore imply that
$`\mathrm{Tr}\left\{\mathrm{exp}[\beta \widehat{\mathrm{\Delta }}]\right\}`$ $`=`$ $`{\displaystyle \frac{1}{|Q|}}{\displaystyle \underset{\eta Q}{}}\mathrm{Tr}\left\{\mathrm{exp}[\beta \widehat{\mathrm{\Delta }}^{(\eta )}]\right\}\mathrm{Tr}\left\{\mathrm{exp}\left[{\displaystyle \frac{\beta }{|Q|}}{\displaystyle \underset{\eta Q}{}}\widehat{\mathrm{\Delta }}^{(\eta )}\right]\right\}`$ (3.54)
$``$ $`e^{\beta d/m}\mathrm{Tr}\left\{\mathrm{exp}[\beta \mathrm{\Delta }]\right\}.`$
We diagonalize $`\mathrm{\Delta }`$ by discrete Fourier transformation on $`^\mathrm{\Lambda }`$. The eigenvalues of $`\mathrm{\Delta }`$ are given by $`\omega (k)`$, where $`k\mathrm{\Lambda }^{}=\frac{2\pi }{L}_L^d`$ is the variable dual to $`x\mathrm{\Lambda }`$. Since $`|\mathrm{\Lambda }^{}|=L^d=|Q|r^d`$, we therefore have
$$\frac{1}{|Q|}\underset{zQ}{}\mathrm{exp}[\beta \mathrm{\Delta }_Q]_{z,z}=\frac{1}{|\mathrm{\Lambda }|}\mathrm{Tr}\left\{\mathrm{exp}[\beta \widehat{\mathrm{\Delta }}]\right\}\frac{e^{\beta d/m}}{|\mathrm{\Lambda }^{}|}\underset{k\mathrm{\Lambda }^{}}{}e^{\beta \omega (k)}.$$
(3.55)
Inequality (3.55) holds for every $`L=r(2m+1)`$, and hence also in the limit $`L\mathrm{}`$. Since the right side of (3.50) is a Riemann sum approximation to the integral in (3.45), this limit yields the asserted estimate (3.45). QED
### 3.4 The Discrete Laplacians on $`\mathrm{\Omega }`$, $`\mathrm{\Omega }^c`$, and their Eigenvalue Sums
In this final subsection, we compare the sum of the eigenvalues of
$$\stackrel{~}{\mathrm{\Delta }}:=P_\mathrm{\Omega }(\mathrm{\Delta })P_\mathrm{\Omega }+P_\mathrm{\Omega }^{}(\mathrm{\Delta })P_\mathrm{\Omega }^{}$$
(3.56)
below $`\mu `$ to the sum of the eigenvalues of $`\mathrm{\Delta }`$ below $`\mu `$, where $`\mathrm{\Omega }\mathrm{\Lambda }`$ is an arbitrary, but henceforth fixed, subset of $`\mathrm{\Lambda }`$, and $`\mathrm{\Omega }^c:=\mathrm{\Lambda }\mathrm{\Omega }`$ is its complement. To this end, we introduce the difference of these eigenvalue sums,
$`\delta E(\mu ,\mathrm{\Omega })`$ $`:=`$ $`\mathrm{Tr}\{[\stackrel{~}{\mathrm{\Delta }}\mu ]_{}\}\mathrm{Tr}\{[\mathrm{\Delta }\mu ]_{}\}`$
$`=`$ $`\mathrm{Tr}\{(\stackrel{~}{\mathrm{\Delta }}\mu )\stackrel{~}{P}_{}\}\mathrm{Tr}\{(\mathrm{\Delta }\mu )P_{}\},`$
where $`\stackrel{~}{P}_{}:=1\mathrm{l}[\stackrel{~}{\mathrm{\Delta }}\mu ]`$ and $`P_{}:=1\mathrm{l}[\mathrm{\Delta }\mu ]`$. We further set $`\stackrel{~}{P}_+:=\stackrel{~}{P}_{}^{}`$ and $`P_+:=P_{}^{}`$. Since $`\stackrel{~}{P}_{}`$ commutes with $`P_\mathrm{\Omega }`$, we have that $`\mathrm{Tr}\{(\stackrel{~}{\mathrm{\Delta }}\mu )\stackrel{~}{P}_{}\}=\mathrm{Tr}\{(\mathrm{\Delta }\mu )\stackrel{~}{P}_{}\}`$, and thus
$`\delta E(\mu ,\mathrm{\Omega })`$ $`=`$ $`\mathrm{Tr}\{(\mathrm{\Delta }\mu )(\stackrel{~}{P}_{}P_{})\}`$
$`=`$ $`\mathrm{Tr}\{[\mathrm{\Delta }\mu ]_{}(\stackrel{~}{P}_{}1\mathrm{l})\}+\mathrm{Tr}\{[\mathrm{\Delta }\mu ]_+\stackrel{~}{P}_{}\}`$
$`=`$ $`\mathrm{Tr}\{[\mathrm{\Delta }+\mu ]_+\stackrel{~}{P}_+\}+\mathrm{Tr}\{[\mathrm{\Delta }\mu ]_+\stackrel{~}{P}_{}\}\mathrm{\hspace{0.33em}0}`$
is manifestly nonnegative. The derivation of a nontrivial lower bound on $`\delta E(\mu ,\mathrm{\Omega })`$ of the form $`\delta E(\mu ,\mathrm{\Omega })\alpha (\mu )|\mathrm{\Omega }|`$, where $`\alpha (\mu )>0`$ is a positive constant which depends only on $`\mu `$ and the spatial dimension $`d1`$ (but not on $`\mathrm{\Omega }`$), is a task that was first addressed by Freericks, Lieb, and Ueltschi in . Shortly thereafter, Goldbaum improved the numerical value for $`\alpha (\mu )>0`$, especially if $`\mu `$ is close to $`2d`$. As a consequence of the estimates in , we have the following lemma.
###### Lemma 3.6 (Freericks, Lieb, and Ueltschi (2002), Goldbaum (2003)).
(i) Let $`\frac{1}{2}<\mathrm{\mu }<4\mathrm{d}`$. There is $`\mathrm{L}_{}(\mathrm{\mu })<\mathrm{}`$ and $`\mathrm{\alpha }(\mathrm{\mu })>0`$ such that, for all $`\mathrm{L}\mathrm{L}_{}(\mathrm{\mu })`$ and all subsets $`\mathrm{\Omega }\mathrm{\Lambda }`$,
$$\delta E(\mu ,\mathrm{\Omega })\alpha (\mu )|\mathrm{\Omega }|.$$
(3.59)
(ii) Let $`0<\mathrm{\mu }\frac{1}{2}`$, and define
$$\alpha (\mu ):=\frac{|S^{d1}|\mu ^{(2+d)/2}}{2^{1+d/2}(2\pi )^d(4d)^5}\text{and}L_{}(\mu ):=\frac{4\pi d}{\mu }.$$
(3.60)
where $`|S^{d1}|`$ is the surface volume of the $`d`$-dimensional sphere. Then, for all $`LL_{}(\mu )`$ and all subsets $`\mathrm{\Omega }\mathrm{\Lambda }=_L^d`$, we have
$$\delta E(\mu ,\mathrm{\Omega })\alpha (\mu )|\mathrm{\Omega }|.$$
(3.61)
*Proof.* We only give the proof of (ii), which amounts to reproducing the proof of Lemma 3.1 in . By $`\{\psi _k\}_{k\mathrm{\Lambda }^{}}^\mathrm{\Lambda }`$ we denote the orthonormal basis (ONB) of eigenvectors of $`\mathrm{\Delta }`$, i.e.,
$$\psi _k(x):=|\mathrm{\Lambda }|^{1/2}e^{ikx},k\mathrm{\Lambda }^{}=\frac{2\pi }{L}_L^d,$$
(3.62)
and we have that $`\mathrm{\Delta }\psi _k=\omega (k)\psi _k`$, with $`\omega (k)=_{\nu =1}^d2\{1\mathrm{cos}(k_\nu )\}`$. Evaluating the traces in Eq. (3.4) by means of this ONB, we obtain
$`\delta E(\mu ,\mathrm{\Omega })`$ $`=`$ $`{\displaystyle \underset{k\mathrm{\Lambda }^{}}{}}\left\{[\mu \omega (k)]_+\psi _k|\stackrel{~}{P}_+\psi _k+[\omega (k)\mu ]_+\psi _k|\stackrel{~}{P}_{}\psi _k\right\}.`$ (3.63)
$``$ $`{\displaystyle \underset{k\mathrm{\Lambda }^{}}{}}[\mu \omega (k)]_+\psi _k|\stackrel{~}{P}_+\psi _k.`$
Let $`\{\phi _j\}_{j=1}^{|\mathrm{\Lambda }|}^\mathrm{\Lambda }`$ be an ONB of eigenvectors of $`\stackrel{~}{\mathrm{\Delta }}`$, i.e., $`\stackrel{~}{\mathrm{\Delta }}\phi _j=e_j\phi _j`$. For any $`k\mathrm{\Lambda }^{}`$ and $`1j|\mathrm{\Lambda }|`$, we observe that
$`\left(e_j\omega (k)\right)^2|\psi _k|\phi _j|^2`$ $`=`$ $`|\psi _k|(\mathrm{\Delta }\stackrel{~}{\mathrm{\Delta }})\phi _j|^2`$ (3.64)
$`=`$ $`|\psi _k|(P_\mathrm{\Omega }\mathrm{\Delta }P_\mathrm{\Omega }^{}+P_\mathrm{\Omega }^{}\mathrm{\Delta }P_\mathrm{\Omega })\phi _j|^2`$
$`=`$ $`|P_\mathrm{\Omega }\mathrm{\Delta }P_\mathrm{\Omega }^{}\psi _k|\phi _j|^2+|P_\mathrm{\Omega }^{}\mathrm{\Delta }P_\mathrm{\Omega }\psi _k|\phi _j|^2`$
$``$ $`|P_\mathrm{\Omega }\mathrm{\Delta }P_\mathrm{\Omega }^{}\psi _k|\phi _j|^2,`$
using that either $`P_\mathrm{\Omega }\phi _j=0`$ or $`P_\mathrm{\Omega }^{}\phi _j=0`$ and that $`P_\mathrm{\Omega }\mathrm{\Delta }P_\mathrm{\Omega }^{}=P_\mathrm{\Omega }\mathrm{\Delta }P_\mathrm{\Omega }^{}`$. Since $`|e_j\omega (k)|4d`$, Eq. (3.64) implies that
$$(4d)^2|\psi _k|\phi _j|^2|b_k|\phi _j|^2,$$
(3.65)
where $`b_k:=P_\mathrm{\Omega }\mathrm{\Delta }P_\mathrm{\Omega }^{}\psi _k`$ is the boundary vector that plays a crucial role in . By summation over all $`j`$ corresponding to eigenvalues $`e_j>\mu `$, we obtain
$$\psi _k|\stackrel{~}{P}_+\psi _k(4d)^2b_k|\stackrel{~}{P}_+b_k,$$
(3.66)
for all $`k\mathrm{\Lambda }^{}`$. Next, the convexity of $`\lambda [\lambda ]_+`$ and the fact that $`\stackrel{~}{P}_+=1\mathrm{l}[\stackrel{~}{\mathrm{\Delta }}>\mu ](4d)^1[\stackrel{~}{\mathrm{\Delta }}\mu ]_+`$ yield
$`b_k|\stackrel{~}{P}_+b_k`$ $``$ $`{\displaystyle \frac{1}{4d}}b_k|[\stackrel{~}{\mathrm{\Delta }}\mu ]_+b_k{\displaystyle \frac{1}{4d}}\left[b_k|(\stackrel{~}{\mathrm{\Delta }}\mu )b_k\right]_+`$ (3.67)
$`=`$ $`{\displaystyle \frac{1}{4d}}\left[b_k|(\mathrm{\Delta }\mu )b_k\right]_+.`$
Now, for any $`x\mathrm{\Omega }`$ there is, by definition, at least one point $`x+e\mathrm{\Omega }^c`$, with $`|e|_1=1`$. Since $`b_k`$ is supported in $`\mathrm{\Omega }`$, we have $`b_k(x+e)=0`$, and thus
$`b_k|(\mathrm{\Delta }\mu )b_k`$ $`=`$ $`{\displaystyle \underset{x\mathrm{\Omega }}{}}\left\{{\displaystyle \underset{|e|_1=1}{}}|b_k(x)b_k(x+e)|^2\mu |b_k(x)|^2\right\}`$ (3.68)
$``$ $`(1\mu ){\displaystyle \underset{x\mathrm{\Omega }}{}}|b_k(x)|^2=(1\mu )b_k^2.`$
Inserting (3.66)โ(3.68) into (3.63), we arrive at
$$\delta E(\mu ,\mathrm{\Omega })\frac{(1\mu )}{(4d)^3}\underset{k\mathrm{\Lambda }^{}}{}[\mu \omega (k)]_+b_k^2.$$
(3.69)
Next, we use that in the sum in (3.69) only those $`k\mathrm{\Lambda }^{}`$ contribute, for which $`\omega (k)=_{\nu =1}^d2\left\{1\mathrm{cos}(k_\nu )\right\}\frac{1}{2}`$, as $`0<\mu 1`$. This implies that $`\mathrm{cos}(k_\nu )\frac{1}{2}`$, for all $`\nu \{1,2,\mathrm{},d\}`$. Hence, for these $`k`$, we have that
$`b_k^2`$ $`=`$ $`{\displaystyle \frac{1}{|\mathrm{\Lambda }|}}{\displaystyle \underset{x\mathrm{\Omega }}{}}\left|{\displaystyle \underset{\sigma =\pm }{}}{\displaystyle \underset{\nu =1}{\overset{d}{}}}e^{i\sigma k_\nu }\mathrm{\hspace{0.25em}1}\mathrm{l}[x+\sigma e_\nu \mathrm{\Omega }^c]\right|^2`$ (3.70)
$``$ $`{\displaystyle \frac{1}{|\mathrm{\Lambda }|}}{\displaystyle \underset{x\mathrm{\Omega }}{}}\left({\displaystyle \underset{\sigma =\pm }{}}{\displaystyle \underset{\nu =1}{\overset{d}{}}}\mathrm{cos}(k_\nu )\mathrm{\hspace{0.25em}1}\mathrm{l}[x+\sigma e_\nu \mathrm{\Omega }^c]\right)^2`$
$``$ $`{\displaystyle \frac{1}{4|\mathrm{\Lambda }|}}{\displaystyle \underset{x\mathrm{\Omega }}{}}1={\displaystyle \frac{|\mathrm{\Omega }|}{4|\mathrm{\Lambda }|}},`$
since there is at least one choice for $`(\sigma ,\nu )`$ such that $`x+\sigma e_\nu \mathrm{\Omega }^c`$. Inserting this estimate into (3.69), we obtain
$$\delta E(\mu ,\mathrm{\Omega })\frac{|\mathrm{\Omega }|}{8(4d)^3}\left(\frac{1}{|\mathrm{\Lambda }^{}|}\underset{k\mathrm{\Lambda }^{}}{}[\mu \omega (k)]_+\right).$$
(3.71)
Now define $`q:๐^d\mathrm{\Lambda }^{}`$ by the preimages
$$q^1(k):=k+[\frac{\pi }{L},\frac{\pi }{L})^d,$$
(3.72)
for $`k\mathrm{\Lambda }^{}`$. In other words, given $`\xi ๐^d`$, the point $`q(\xi )\mathrm{\Lambda }^{}`$ is the closest point to $`\xi `$. In particular, $`|\xi q(\xi )|_{\mathrm{}}\frac{\pi }{L}`$, which implies that $`|\omega (q(\xi ))\omega (\xi )|\frac{2\pi d}{L}`$, by Taylorโs theorem. Hence,
$`{\displaystyle \frac{1}{|\mathrm{\Lambda }^{}|}}{\displaystyle \underset{k\mathrm{\Lambda }^{}}{}}[\mu \omega (k)]_+`$ $`=`$ $`{\displaystyle _{๐^d}}\left[\mu \omega (q(\xi ))\right]_+{\displaystyle \frac{d^d\xi }{(2\pi )^d}}`$ (3.73)
$``$ $`{\displaystyle _{๐^d}}\left[\mu 2\pi dL^1\omega (\xi )\right]_+{\displaystyle \frac{d^d\xi }{(2\pi )^d}}.`$
Since, by assumption, $`\frac{2\pi d}{L}\frac{2\pi d}{L_{}}=\frac{\mu }{2}`$ and $`\omega (\xi )\xi ^2`$, we have
$$_{๐^d}\left[\mu 2\pi dL^1\omega (\xi )\right]_+d^d\xi _{๐^d}\left[\frac{\mu }{2}\xi ^2\right]_+d^d\xi =\frac{|S^{d1}|}{2^{d/2}d(d+2)}\mu ^{1+(d/2)}.$$
(3.74)
Inserting (3.73)โ(3.74) into (3.71), we arrive at the asserted estimate. QED
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# Dependence of the 12C(๐พโ,pd) reaction on photon linear polarisation
Present address\] Edinburgh University, UK.
## Abstract
The sensitivity of the <sup>12</sup>C$`(\stackrel{}{\gamma },pd)`$ reaction to photon linear polarisation has been determined at MAMI, giving the first measurement of the reaction for a nucleus heavier than <sup>3</sup>He. Photon asymmetries and cross sections were measured for $`E_\gamma `$=170 to 350 MeV. For $`E_\gamma `$ below the $`\mathrm{\Delta }`$ resonance, reactions leaving the residual <sup>9</sup>Be near its ground state show a positive asymmetry of up to 0.3, similar to that observed for <sup>3</sup>He suggesting a similar reaction mechanism for the two nuclei.
Three-body forces have consequences in many fields of physics. The study of photon induced proton-deuteron knockout from nuclei may give valuable information on the three-body interaction in the nucleus, since the direct mechanisms which contribute may be related to those thought to be involved in the three-nucleon forceSkibinski et al. (2003a, b); Deltuva et al. (2004); Laget (1988). However, as well as the direct 3-nucleon process ($`3N`$) there will be contributions from initial photon absorption by a single nucleon ($`1N`$), two-nucleons ($`2N`$) and two-step $`3N`$ processes such as initial real pion production on one nucleon followed by reabsorption by a nucleon pair. Clearly, to extract reliable information from $`(\gamma ,pd)`$ measurements, the relative contributions from each of these mechanisms should be well understood.
The ($`\gamma `$,pd) reaction has received significant theoretical interest in recent years, mainly motivated by the possibility of obtaining information on the nature of the 3-nucleon force (3NF). Detailed <sup>3</sup>He calculations based on exact solutions of the three-particle scattering equations in the initial and final states have been carried out for photon energies up to 140 MeVSkibinski et al. (2003a); Deltuva et al. (2004). These show that the inclusion of a 3NF has a large effect on the magnitude of the cross section, increasing the predictions by up to a factor of two at the top end of this $`E_\gamma `$ range. A microscopic theoretical treatment of the <sup>3</sup>He($`\gamma `$,pd) reaction, which includes contributions from $`1N`$, $`2N`$ and $`3N`$ mechanisms, has been developed by LagetLaget (1988). The $`3N`$ mechanisms include contributions from virtual and real pion exchange. Lagetโs treatment relies on a factorisation approximation to simplify the computation but is applicable up to higher photon energies.
On the experimental side, most measurements of the ($`\gamma `$,pd) reaction have been made using <sup>3</sup>He targetsIsbert et al. (1994); Sober et al. (1983); Argan et al. (1975); Gassen et al. (1981); Picozza et al. (1970); Kolb et al. (1994). An important feature of the measured excitation functions is that they show no evidence of structure for photon energies in the $`\mathrm{\Delta }`$(1232) resonance region. Also, the centre-of-momentum (CM) proton angle distributions are forward peaked and fall off rapidly with increasing angle up to $``$70 with a flatter distribution at more backward angles. The features are moderately well described by the Laget model when $`1N`$, $`2N`$ and two-step $`3N`$ (including only real $`\pi `$ exchange) mechanisms are includedLaget (1988). Laget notes that his model accounts less well for the <sup>3</sup>He$`(\gamma ,pd)`$ data than it does for $`\pi `$ induced processes involving the A=3 nuclei and suggests two additional photon couplings, both involving two highly virtual mesons, which could be responsible. Above 100 MeV the $`2N`$ and two-step $`3N`$ mechanisms are predicted to dominate with the $`2N`$ mechanisms only giving large contributions to the cross section at forward CM proton anglesLaget (1988); Isbert et al. (1994). Above $``$150 MeV the two-step 3N mechanism is predicted to provide most of the cross section for CM proton angles backwards of $``$70.
Studies of the $`(\gamma ,pd)`$ reaction for A$`>`$3 targets have been carried out only on <sup>16</sup>OHartmann et al. (1973) and <sup>12</sup>CMcAllister et al. (1999). Both measurements show a photon energy dependence similar to that observed in <sup>3</sup>He with no prominent enhancement for $`E_\gamma `$ around the $`\mathrm{\Delta }`$ resonance. This behaviour is in contrast to photon induced $`pp`$, $`pn`$, $`p\pi `$ and $`ppn`$MacGregor et al. (1998); Branford et al. (2000); Audit et al. (1997); Watts et al. (2003) knockout reactions where the $`\mathrm{\Delta }`$ plays a prominent role. The <sup>12</sup>C($`\gamma ,pd)`$ missing energy spectra obtained in Ref. McAllister et al. (1999) exhibit significant strength close to the reaction threshold, and the recoil momentum spectra of the (Aโ3) nucleus at low missing energies are consistent with those predicted if it were a spectator to the knockout of three 1p-shell nucleons.
The photon asymmetry for the $`(\stackrel{}{\gamma },pd)`$ reaction has only been measured previously for <sup>3</sup>HeBelyaev et al. (1984); Fabbri et al. (1972). These measurements showed a positive asymmetry which ranged from around 0.2 to 0.5 over the sampled photon energy region of 90-350 MeV and CM proton angle range of 60-135. Comparison with a simple Faddeev calculationBelyaev et al. (1984) which neglects meson exchange currents, $`\mathrm{\Delta }`$ contributions and final state interactions gave limited agreement with the experimental data.
The three nucleon photoabsorption mechanisms which operate in the $`(\gamma ,pd)`$ reaction also contribute in the $`(\gamma ,ppn)`$ reaction but with less restrictive spin and isospin conditions for the final state particlesCarrasco et al. (1994); Audit et al. (1997). Recent Faddeev theoretical calculations for <sup>3</sup>HeSkibinski et al. (2003a) indicate that the $`(\gamma ,ppn)`$ reaction is also sensitive to the nature of the 3NF. Several measurements of this reaction have been made in the last decadeRuth. et al. (1994); Audit et al. (1997); Watts et al. (2003); Sarty et al. (1993). Polarised photon measurements of the <sup>3</sup>He$`(\stackrel{}{\gamma },p)X`$ reactionRuth. et al. (1994) for regions of nucleon momenta where the two-step 3N processes were expected to dominate showed a negative asymmetry of up to 0.2 which is not well described by the Laget model. A recent comparison of the <sup>12</sup>C$`(\gamma ,ppn)`$ reactionWatts et al. (2003), with model calculationsCarrasco et al. (1994) in restricted kinematics for the detected nucleons gave clear evidence for the existence of a direct $`3N`$ mode.
The present work is the first measurement of the <sup>12</sup>C($`\stackrel{}{\gamma }`$,pd) reaction. The sensitivity of the cross section to photon linear polarisation is expected to give valuable constraints on the reaction mechanisms for different photon energies and in different excitation energy regions of the (Aโ3) nucleus.
The experiment was carried out at the 855 MeV Mainz microtron (MAMI-B) Herminghaus. (1990); Walcher (1990) using the Glasgow tagged-photon spectrometerAnthony et al. (1985); Hall et al. (1996) in conjunction with two plastic scintillator arrays, PiP and TOFMacGregor et al. (1996); Grabmayr et al. (1998), set up as described in Ref. Powrie et al. (2001); Franczuk et al. (1999). Polarised photons were produced by coherent bremsstrahlung in a thin diamond radiatorLohmann et al. (1994); Kraus et al. (1997); Rambo et al. (1998); Natter et al. (2002). The polarisation orientation of the photons was flipped between horizontal and vertical every few minutes. Three angular settings of the diamond were used for which the main coherent peak covered the photon energy ranges 170โ220, 220โ280 and 300โ350 MeV with corresponding average linear polarisation ($`P`$) of 59.5%, 49.5% and 42.5% respectively.
Protons with kinetic energies 31-270 MeV were detected in the charged particle hodoscope PiP covering the polar angular range $`\theta `$=51-129 and azimuthal angular range of $`\varphi `$=$`\pm `$23. Coincident deuterons leaving the target with energies above $``$45 MeV were detected in TOF which determined particle energies by time-of-flight. The TOF detectors covered $`\theta `$=10.0-175.0. The deuterons were separated from other charged hadrons in TOF by selecting events from a 2-D plot of inverse speed versus pulse height, as described in Ref. McAllister et al. (1999). The number of random deuteron coincidences in TOF was found to be negligible. The average measured missing energy resolution extracted from D$`(\gamma ,pn)`$ was found to be $``$7 MeV. For $`(\gamma ,pd)`$ the average resolution is better due to the slower flight times for the heavier particles and is estimated to be $``$5 MeV. The $``$3% background of events not originating from reactions in the target was measured in runs with the target removed. The total systematic error in the asymmetry is estimated to be $`\mathrm{\Delta }\mathrm{\Sigma }`$ = $`\pm `$0.05$`\mathrm{\Sigma }`$Powrie et al. (2001). The systematic uncertainties in the measured cross sections are estimated to be up to $`\pm `$8%McAllister et al. (1999).
For each goniometer setting all events due to photons in the coherent peak region were used to produce average $`(\stackrel{}{\gamma },pd)`$ cross sections ($`\sigma `$) and asymmetries ($`\mathrm{\Sigma }`$) for the events which were within the geometrical and energy acceptances of the PiP-TOF detector systems. The asymmetry is defined as $`\mathrm{\Sigma }=\frac{1}{P}\frac{\sigma ^{}\sigma ^{}}{\sigma ^{}+\sigma }`$ where $`\sigma ^{}`$ ($`\sigma ^{}`$) is the measured cross section for reactions in which the horizontal detector plane is parallel(perpendicular) to the electric vector of the polarised photons. The reduction of the asymmetry arising from the particle detectors having a finite $`\varphi `$ acceptance around the horizontal plane is estimated to be $``$10%. A further reduction which arises from the smearing of the photon polarisation direction in the (3N+$`\gamma `$) CM frame due to the initial 3N Fermi momentum is estimated to be $``$5%. The magnitude of the presented asymmetries have been increased by 15% to account for these effects.
The cross section as a function of $`E_m`$ is shown in Fig. 1 for three different $`E_\gamma `$ bins. Missing energy is defined as $`E_m=E_\gamma T_1T_2T_r`$ where $`E_\gamma `$ is the incident photon energy, $`T_1`$ and $`T_2`$ are the energies of the two detected particles and $`T_r`$ is the (typically small) energy of the recoiling system which is calculated from its momentum $`๐_r=๐_\gamma ๐_p๐_d`$. The Q-value for the <sup>12</sup>C$`(\stackrel{}{\gamma },pd)`$ reaction is โ31.7 MeV. The missing energy spectra show strength near to the Q-value and also a peak at $``$45 MeV. Above this peak, the $`(\gamma ,pd)`$ missing energy distribution rises to a second maximum, which becomes higher and wider as $`E_\gamma `$ increases. Similar general features in the missing energy are seen in the earlier work of Ref. McAllister et al. (1999), but with weaker indications of the peak at $``$45 MeV, probably due to its poorer statistical accuracy and inferior resolution. A $`(\gamma ,pd)`$ reaction which involves the knockout of three (1p) shell nucleons leaving a residual (Aโ3) spectator nucleus would populate missing energies up to $``$50 MeV. The missing energies expected if the residual <sup>9</sup>Be is left in one of its four low lying and relatively long lived states are also indicated in Fig. 1. The peak at $``$45 MeV is seen to occur at energies close to the lowest lying T=$`\frac{3}{2}`$ states at 14.4 and 17.0 MeV. Also shown in Fig. 1 are the scaled <sup>12</sup>C$`(\gamma ,pp)`$ and <sup>12</sup>C$`(\gamma ,pn)`$ cross sections obtained with the same detector setup as the present measurement (this data was already analysed in Ref. Powrie et al. (2001)). Above $`E_m`$60 MeV the $`(\gamma ,NN)`$ reactions scale with $`E_m`$ and $`E_\gamma `$ in a similar way to the present <sup>12</sup>C$`(\gamma ,pd)`$ data. At low $`E_m`$ the channels show different behaviour. The $`(\gamma ,pp)`$ data falls off more rapidly than $`(\gamma ,pd)`$ and the $`(\gamma ,pn)`$ cross section has a large peak at $`E_m30`$MeV (off scale in Fig. 1) due to a large direct two-nucleon knockout contribution.
The photon asymmetry for the <sup>12</sup>C$`(\gamma ,pd)`$ reaction is presented in Fig. 2. The $`E_m<`$40 MeV region emphasizes $`(\gamma ,pd)`$ reactions leading to the ground state and low lying excited states of <sup>9</sup>Be, which are all T=$`\frac{1}{2}`$. A positive asymmetry is observed for photon energies up to $`E_\gamma `$280 MeV, while at higher photon energies in the $`\mathrm{\Delta }`$ resonance region the asymmetry becomes negative. The asymmetry for <sup>3</sup>He($`\stackrel{}{\gamma },pd`$) for $`\theta _p`$=90 and 110 in the CM frameBelyaev et al. (1984) (both corresponding to lab proton angles covered by the PiP detector in the present measurement) are also shown in Fig. 2(a). The magnitude and sign of the <sup>3</sup>He asymmetry is similar to the $`E_m<`$40 MeV <sup>12</sup>C data for photon energies up to $``$270 MeV. At higher $`E_\gamma `$ the <sup>3</sup>He data do not show the negative or small asymmetries indicated in the <sup>12</sup>C($`\gamma ,pd`$) data.
The $`E_m`$=40-50 MeV cut (Fig. 2(b)) emphasizes $`(1p)^3`$ knockout events leading to higher excited states and includes the peak region visible in the cross section as a function of missing energy (Fig. 1). The asymmetry in this region is generally negative or small below $``$300 MeV in contrast to the positive asymmetry observed at lower missing energies for both <sup>12</sup>C and <sup>3</sup>He($`\stackrel{}{\gamma },pd)`$. The asymmetry for the $`E_m`$=50-100 MeV region shows the same general trends as the $`E_m`$=40-50 MeV data, albeit smaller in magnitude. The asymmetries for both these higher $`E_m`$ regions can be seen to show features which are very similar to those observed for the <sup>12</sup>C($`\stackrel{}{\gamma },NN`$) reactions at high missing energyPowrie et al. (2001), suggestive of similar underlying reaction mechanisms. The model of Ref. Carrasco et al. (1994) explains the $`(\gamma ,NN)`$ cross section in this missing energy region largely as the result of detecting two of the three (or more) nucleons produced by a two-step 3N process or by initial photon absorption on a two-nucleon pair followed by final state interactionsLamparter et al. (1996); Watts et al. (2000). The same processes, where one of the outgoing nucleons picks up an additional nucleon from the residual nucleus, can probably explain the similar $`E_m`$ distribution and asymmetry of the <sup>12</sup>C$`(\gamma ,pd)`$ reaction. As these mechanisms involve more than three nucleons they have no analogue in the reaction on <sup>3</sup>He.
A comparison of the $`(\gamma ,pp)`$ and $`(\gamma ,pd)`$ $`E_m`$ distributions below 60 MeV may suggest the pickup process discussed in the previous paragraph still provides a significant background contribution in this region. However, this extrapolation should be considered an upper limit as the dominant mechanism of the $`(\gamma ,pp)`$ reaction changes to direct two-proton emission following photon absorption on two-nucleonsLamparter et al. (1996); Watts et al. (2000) at low $`E_m`$, which only has significant subsequent pickup probability at more forward proton angles than sampled hereLaget (1988). This is supported by the comparison in Fig. 2 of the <sup>12</sup>C$`(\stackrel{}{\gamma },pd)`$ asymmetry with the asymmetry of the <sup>12</sup>C($`\stackrel{}{\gamma },pp`$), and <sup>12</sup>C($`\stackrel{}{\gamma },pn`$) reactions for $`E_m`$40 MeVPowrie et al. (2001). Both the reactions show a negative asymmetry. The positive asymmetries observed in <sup>12</sup>C($`\stackrel{}{\gamma },pd`$) therefore argue against a large direct feeding of strength from ($`\gamma ,pp`$) and ($`\gamma ,pn`$) reactions through subsequent pickup reactions at the lowest missing energies.
The selection of <sup>12</sup>C$`(\stackrel{}{\gamma },pd)`$ events with low missing energy should enhance the contribution of processes which involve only the three detected nucleons while the (Aโ3) nucleus spectates. The similar asymmetry observed for $`(\gamma ,pd)`$ reaction in <sup>12</sup>C in this $`E_m`$ region and <sup>3</sup>He (Fig. 2(a)) suggests similar reaction mechanisms in both nuclei. The Laget model predicts that the two-step 3N mechanism involving the initial production of an on-shell pion is the dominant mechanism for the $`E_\gamma `$ and $`\theta _p`$ sampled in the present experiment. The published Laget calculations for <sup>3</sup>HeLaget (1988) do not present results for the ($`\gamma ,pd`$) asymmetry, but an indication of its behaviour is sought here by examining the asymmetry in the initial p($`\gamma ,N`$)$`\pi `$ stage of the dominant 3N two-step mechanism. Since the $`(\gamma ,pd)`$ process involves deuteron formation from the recoiling nucleon in this process, then the underlying ($`\gamma ,N`$)$`\pi `$ asymmetry should be reflected in the $`(\gamma ,pd)`$ data. This feeding of the underlying ($`\gamma ,N`$)$`\pi `$ asymmetry to the final state particles in the two-step 3N mechanism was already indicated in the $`(\gamma ,NN)`$ measurementsPowrie et al. (2001) at high missing energy.
To determine whether these considerations apply to our data we used the predictions for the $`\stackrel{}{\gamma }NN\pi `$ asymmetry obtained using the MAID codeDreschsel et al. (1999), which is based on a unitary isobar parameterisation of experimental data and accounts for the non-resonant and resonant parts of the pion photoproduction amplitude. Both p($`\stackrel{}{\gamma },\pi ^+)`$n and p($`\stackrel{}{\gamma },\pi ^0)`$p were calculated at a pion CM breakup angle of 55, which for photon energies above 200 MeV results in recoiling nucleon angles in regions where the deuteron yield in the present <sup>12</sup>C$`(\stackrel{}{\gamma },pd)`$ data is largest ($``$40-60). The predictions are presented in Fig. 2(a).
Neither of the full p($`\gamma ,\pi `$)$`N`$ calculations can give a simple explanation of the low missing energy <sup>12</sup>C$`(\stackrel{}{\gamma },pd)`$ asymmetry in terms of the two-step 3N process. As the $`\mathrm{\Delta }`$ contribution to the initial pion production vertex for the two-step 3N process is suppressed due to an isospin restrictionLaget (1988), MAID calculations with the $`\mathrm{\Delta }`$ contribution removed are also shown in Fig. 2. While the p($`\gamma ,\pi ^0`$)p MAID prediction with no $`\mathrm{\Delta }`$ contribution comes closer to the $`(\stackrel{}{\gamma },pd)`$ asymmetry, the MAID cross sections suggest a dominance for the p($`\gamma ,\pi ^+`$)n process for which the asymmetry is negative. It therefore seems unlikely that there is a simple explanation of the <sup>12</sup>C$`(\stackrel{}{\gamma },pd)`$ or <sup>3</sup>He$`(\stackrel{}{\gamma },pd)`$ asymmetries in terms of the two-step 3N mechanism.
It is clear that theoretical input is required to fully account for the various mechanisms and their possible interference. Of particular interest would be estimates which include 3N mechanisms involving heavier mesons and also the two-meson couplings suggested by Laget. Such processes produce a shorter range interaction than the two-step 3N mechanisms and are thought to be important in the 3N interaction.
In summary this first determination of the <sup>12</sup>C($`\stackrel{}{\gamma },pd`$) asymmetry shows that reactions leading to low lying states in <sup>9</sup>Be proceed through the photon interacting with the detected nucleons in a similar manner to the <sup>3</sup>He($`\stackrel{}{\gamma },pd`$) reaction for $`E_\gamma `$ below the $`\mathrm{\Delta }`$ resonance. The asymmetries at higher missing energy do not resemble <sup>3</sup>He($`\stackrel{}{\gamma },pd`$) and have a plausible explanation in terms of multistep processes involving more than three nucleons. These new results will provide valuable constraints on the reaction mechanisms for $`(\gamma ,pd)`$ in heavier nuclei and their potential to be used in learning about the correlated behaviour of three nucleons in a nucleus.
###### Acknowledgements.
This work was supported by the UK EPSRC, the British Council, the DFG (Mu 705/3, SSP 1043), BMFT (06 Tuยจ 656), DAAD (313-ARC-IX-95/41), the EC \[SCI.0910.C(JR)\] and NATO (CRG 970268).
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# An hybrid Tykhonov method for neutron spectrum unfolding
## 1 Introduction
Following the ideas of A.N. Tykhonov about ill-posed problems inversion, we extend the iterative regularization method by including alternatively between two iterations least square gradient looking terms, and with a modification of the parametrization of the regularizing term. The merits of this approach are a fast convergence to a satisfying solution, stable and unsensitive to the noise fluctuations, and a complete independence with respect to the initial guess. With this method, we unfold one dimensional data set of measurements with a Bonner Sphere system to obtain a wide neutron spectrum estimation. Two experimental cases are presented : a complex thermal neutron wide spectrum and a single monoernegetic 565 keV neutron spectrum. The evolution of the criteria parameter is discussed and also the systematics and statistics error bars calculations.
## 2 Theoretical approach
Once we have a set of 11 measurements, by spline interpolation, we obtain a vector $`y`$ of dimension 64 $`(y^{+64})`$; this is justified by the smoothness of the data set and by the convolution that erodes most of the details of the initial spectrum. The measurement function $`yC^{\mathrm{}}`$ and is very smooth. The spline interpolation produces a 64 bin vector, with continous derivatives at each of the 11 knots and local second degree polynomial interpolation. These 11 measurements corresponds to the responses to a neutron spectrum of an Helium<sup>3</sup> neutron detector with various polyethylene spheres, that are equivalent to different filters with response matrix A. With $`x^{+64}`$ as the unknown neutron spectrum, $`y`$ as the measurements, the convolution equation of the system stands as :
$`y=Ax`$ (1)
$`x^{+64}y^{+64}A^{+64}\times ^{+64}`$
As the problem is ill-posed, neither the direct inversion, nor the least square solution are satisfying solutions, because they are irregular and highly sensitive to the fluctuations of the measurementโs uncertainties. That means that continuous variations of $`y`$ does not lead to continuous variations in $`x`$.
Direct inversion :
$`\text{ }\text{estimator : }\stackrel{~}{x}=A^1y`$ (2)
Least square solution :
$`\text{criterion : }Q=|Axy|^2`$ (3)
$`\text{estimator : }\stackrel{~}{x}=[A^tA]^1A^ty`$ (4)
$`(A^t\text{: Transposed matrix})`$
### 2.1 Iterative Tykhonov approach
A better solution can be achieved through the scheme of the basic Tykhonov method , which is more regular and less sensitive to noisy fluctuations. A regularization term is added to the least square criterion through a Lagrange multiplicator method. This term ensures the continuity of the solution and its closedness to an initial guess $`x_0`$, which can be a โnot so badโ first solution.
$`Q_\lambda =|Axy|^2+\lambda |xx_0|^2\lambda [0,+\mathrm{}]`$ (5)
$`\stackrel{~}{x}=[A^tA+\lambda I]^1[A^ty+\lambda x_0]`$ (6)
The demonstration of the calculation of the estimator (equation 6) can be found as follows :
Assume $`\stackrel{~}{x}`$ to be the global minimum of the criterion $`Q_\lambda `$.
$`x^{+64}:Q_{\lambda ,\stackrel{~}{x}}Q_{\lambda ,x}`$ (7)
This is true, moreover in the case of small variations around $`\stackrel{~}{x}`$ :
$`x=\stackrel{~}{x}+\alpha hh^{+64}\alpha `$
$`Q_{\lambda ,\stackrel{~}{x}}Q_{\lambda ,\alpha ,h}`$
$`Q_{\lambda ,\alpha ,h}Q_{\lambda ,\stackrel{~}{x}}0`$
$`|A(\stackrel{~}{x}+\alpha h)y|^2+\lambda |\stackrel{~}{x}+\alpha hx_0|^2|A\stackrel{~}{x}y|^2\lambda |\stackrel{~}{x}x_0|^20`$
Due to the definition of the scalar product :
$`\alpha ^{}(Ah)^t[A\stackrel{~}{x}y]+\alpha [A\stackrel{~}{x}y]^tAh+|\alpha |^2|Ah|^2+\alpha ^{}\lambda h^t[\stackrel{~}{x}x_0]+`$
$`\alpha \lambda [\stackrel{~}{x}x_0]^th+\lambda |\alpha |^2|h|^20`$
$`(\alpha ^{}h^t)(A^t[A\stackrel{~}{x}y]+\lambda [\stackrel{~}{x}x_0])+\alpha ([A\stackrel{~}{x}y]^tA+\lambda [\stackrel{~}{x}x_0]^t)h0`$
$`\text{That is independent of }h^t\text{ and }\alpha \text{ ; then it implies :}`$
$`(A^tA\stackrel{~}{x}A^ty+\lambda \stackrel{~}{x}\lambda x_0)=0`$
$`[A^tA+\lambda I]\stackrel{~}{x}[A^ty+\lambda x_0]=0`$
$`\stackrel{~}{x}=[A^tA+\lambda I]^1[A^ty+\lambda x_0]`$
The uniqueness and existence of a bounded solution is also proved in Bertero . The Tykhonov iterative approach is achieved by successive refinement of the solution with several iterations, beginning with an initial guess $`x_0`$, the solution $`x_1`$ satisfies better the criterion of Tykhonov than $`x_0`$, and so on :
$`Q_{\lambda ,n}=|Ax_ny|^2+\lambda |x_nx_{n1}|^2\lambda [0,+\mathrm{}]`$ (8)
$`\stackrel{~}{x}_n=[A^tA+\lambda I]^1[A^ty+\lambda x_{n1}]`$ (9)
$`Q_{\lambda ,1}Q_{\lambda ,2}\mathrm{}Q_{\lambda ,n}Q_{\lambda ,n+1}\mathrm{}`$
### 2.2 Modification of the running parameter $`\lambda `$
In order to help the computer calculation, a slight modification is done with the modification of the running interval of the parameter $`\lambda `$. We also decided that the two terms of the $`Q_{\lambda ,n}`$ criterion have to be of equal importance :
$`\lambda [0,+\mathrm{}]\lambda [0,1]`$ (10)
$`Q_{\lambda ,n}=(1\lambda )|Ax_ny|^2+\lambda |x_nx_{n1}|^2\lambda [0,1]`$ (11)
$`\stackrel{~}{x}_n=[(1\lambda )A^tA+\lambda I]^1[(1\lambda )A^ty+\lambda x_{n1}]`$ (12)
Letโs have $`x_0`$ as an initial guess, and lets compute $`x_1`$. The case $`\lambda =0`$ gives :
$`Q_{0,1}=|Ax_1y|^2`$ (13)
$`\stackrel{~}{x}_{0,1}=[A^tA]^1A^ty\text{ }\text{(least square solution)}`$ (14)
$`0Q_{0,1}`$ (15)
On the other side, the case $`\lambda =1`$ gives : $`Q_{1,1}=|x_1x_0|^2`$ with the solution $`x_1=x_0`$ then $`Q_{1,1}=0`$ then $`\lambda =1`$ is a minimorum and all the iterations stand for an identity operator, then $`n^{}x_n=x_0`$ independently of $`x_0`$ ! So this method gives always a wrong solution, because all initial guesses are available solutions without any restriction. We fixed this problem by modifying the criteria in the following way :
$`Q_{\lambda ,n}={\displaystyle \frac{(1\lambda )|Ax_ny|^2+\lambda |x_nx_{n1}|^2}{(1\lambda )\lambda }}\lambda ]0,1[`$ (16)
$`Q_{\lambda ,n}={\displaystyle \frac{|Ax_ny|^2}{\lambda }}+{\displaystyle \frac{|x_nx_{n1}|^2}{1\lambda }}`$ (17)
the estimator is unchanged :
$$\stackrel{~}{x}_n=[(1\lambda )A^tA+\lambda I]^1[(1\lambda )A^ty+\lambda x_{n1}]$$
(18)
but $`Q_{0,n}=+\mathrm{}`$ and $`Q_{1,n}=+\mathrm{}`$
Also we can note : $`\lambda ]0,1[`$ , $`Q_{\lambda ,n}<+\mathrm{}`$ $``$ an actual minimorum of $`Q_{\lambda ,n}`$ and the false minimorum $`Q_{1,n}`$ disappears!
### 2.3 Possible generalization to derivatives of order m
Basically, the Tykhonov method can be generalized to derivatives of order m of $`x_n`$ instead of using the simple squared modulus . These derivatives $`D^{(m)}x_n`$ are the derivatives of order m of the spectrum $`x_n(E)`$ regarding the energy binning $`E`$ : $`D^{(m)}x_n=\frac{d^{(m)}}{dE}x_n(E)`$
This generalization to greater orders of derivatives is given by :
$`Q_{\lambda ,n,m}={\displaystyle \frac{|Ax_ny|^2}{\lambda }}+{\displaystyle \frac{|D^{(m)}x_nD^{(m)}x_{n1}|^2}{1\lambda }}`$ (19)
$`\stackrel{~}{x}_n=[(1\lambda )A^tA+\lambda D^{(m)t}D^{(m)}]^1[(1\lambda )A^ty+\lambda D^{(m)t}D^{(m)}x_{n1}]`$ (20)
But we have found that the convergence to a stable solution is rather better and faster with a simple function without the $`m^{th}`$ derivative, so we decided to restrict our calculation only to order 0.
### 2.4 Hybridation with the least square method
Three main concerns guided the development of this unfolding iterative method : a fast converging method, a better solution than the classical Tykhonov method or than the least square method and the last one, a complete independence of the solution regards to the initial guess. The iterative Tykhonov approach ensures the last concern but not the first two. So we decided to introduce some accelerating (focusing) algorithm between two classical Tykhonov iterations. Let suppose that $`x_{n1}`$ and $`x_n`$ are two successive solutions obtained with this Tykhonov algorithm, we create some kind of derivative term regarding the iteration parameter n:
$`x_n/n=x_nx_{n1}`$
$`\underset{n+\mathrm{}}{lim}{\displaystyle \frac{x_n}{n}}=0x^{+64};\underset{n+\mathrm{}}{lim}x_n=x`$
$`\text{Then we minimize }Q_{\mu ,n+1}`$
$`Q_{\mu ,n+1}=|Ax_{n+1}y|^2\text{ with }x_{n+1}=x_n+\mu {\displaystyle \frac{x_n}{n}};\mu `$
$`Q_{\mu ,n+1}=|A(x_n+\mu {\displaystyle \frac{x_n}{n}})y|^2`$ (21)
So after every two Tykhonov unfoldings, we compute one solution closed to the least square solution.
At the end of the convergence, the reconstructed spectrum $`y_n=Ax_n`$ is closer to the measured spectrum $`y`$ than with use of the classical Tykhonov iterative scheme, which presents the compromise of an enhancement of the stability of the solution toward the noise fluctuations while the price to pay is a loss of the accuracy and an excursion from the measured spectrum $`y`$. By the way, this approach converges faster than a classical iterative Tykhonov method. This hybridation algorithm acts as a focalisation method with respect to the requested measurement $`y`$, and also as a great accelerator of convergence.
### 2.5 Iteration schema
The iterative schema can be described in the following way. Letโs have $`x_0`$ as any initial guess.
iteration 1 : (22)
$`Q_{\lambda ,1}={\displaystyle \frac{|Ax_1y|^2}{\lambda }}+{\displaystyle \frac{|x_1x_0|^2}{(1\lambda )}}`$
$`\stackrel{~}{x}_1=[(1\lambda )A^tA+\lambda I]^1[(1\lambda )A^ty+\lambda x_0]`$ (23)
iteration 2 : (24)
$`Q_{\lambda ,2}={\displaystyle \frac{|Ax_2y|^2}{\lambda }}+{\displaystyle \frac{|x_2x_1|^2}{(1\lambda )}}`$
$`\stackrel{~}{x}_2=[(1\lambda )A^tA+\lambda I]^1[(1\lambda )A^ty+\lambda x_1]`$ (25)
iteration 3: (26)
$`Q_{\mu ,3}=|A(x_2+\mu {\displaystyle \frac{x_2}{n}})y|^2`$
$`\stackrel{~}{x}_3=x_2+\mu {\displaystyle \frac{x_2}{n}}`$ (27)
iteration 4 : (28)
$`Q_{\lambda ,4}={\displaystyle \frac{|Ax_4y|^2}{\lambda }}+{\displaystyle \frac{|x_4x_3|^2}{1\lambda }}`$
$`\stackrel{~}{x}_4=[(1\lambda )A^tA+\lambda I]^1[(1\lambda )A^ty+\lambda x_3]`$
$`\mathrm{}`$
iteration n : $`n1\text{ modulo }3`$ (30)
$`Q_{\lambda ,n}={\displaystyle \frac{|Ax_ny|^2}{\lambda }}+{\displaystyle \frac{|x_nx_{n1}|^2}{1\lambda }}`$
$`\stackrel{~}{x}_n=[(1\lambda )A^tA+\lambda I]^1[(1\lambda )A^ty+\lambda x_{n1}]`$ (31)
iteration n+1 : $`n+12\text{ modulo }3`$ (32)
$`Q_{\lambda ,n+1}={\displaystyle \frac{|Ax_{n+1}y|^2}{\lambda }}+{\displaystyle \frac{|x_{n+1}x_n|^2}{1\lambda }}`$
$`\stackrel{~}{x}_{n+1}=[(1\lambda )A^tA+\lambda I]^1[(1\lambda )A^ty+\lambda x_n]`$ (33)
iteration n+2 : $`n+20\text{ modulo }3`$ (34)
$`Q_{\mu ,n+2}=|A(x_{n+1}+\mu {\displaystyle \frac{x_{n+1}}{n}})y|^2`$
$`\stackrel{~}{x}_{n+2}=x_{n+1}+\mu {\displaystyle \frac{x_{n+1}}{n}}`$
$`\mathrm{}`$
We have seen with the help of a computer that in all cases, this iterative schema converges to one unique solution, but we have not build the mathematical proof of this convergency! Practically, on a DEC-Compaq Alpha Workstation working with double precision 64 bits numbers, the resulting estimators are always stable after 300 iterations ( 1 minute of computing time) and oftenly around 80 iterations, it is already stable. That is why we decided to put a cut on the maximum number of iterations at 300. We also decided to restrict the binning (in energy and interpolated sphere) over 64 bins, because working with 32 or 128 bins gives either a too bad energy resolution or a too slow computation, thus 64 seems to be an optimum.
### 2.6 Positivity of the solution
Another important aspect of the unfolding is treated without any subtility: positivity . Of course each bin of the spectrum must have a positive content, on the other case the spectrum would be unphysical. The positivity is introduced abruptly at each step of the iteration with the following non-linearity:
$$\stackrel{~}{x}_{n,i}Sup[0,\stackrel{~}{x}_{n,i}]$$
(36)
### 2.7 Calculation of the statistics and systematics fluctuations
The statistics fluctuations are estimated with a Monte-Carlo method. For one set of 11 measurements with the Bonner Sphere system,11 uncertainties of these measurements are calculated taking into account the number of observed counts and the duration of the measures. The propagation through the unfolding process of the statistics fluctuations of the measurements is computed with 100 sets of 64 data $`y`$, which are interpolated with the spline method from these 11 knots. The random generation of the 11 knots is performed with a gaussian law centered onto the 11 actual measurements and with a standard deviation corresponding to the error bars directly measured with the instrument. The unfolding of the 100 data sets $`y`$ gives 100 $`\stackrel{~}{x}`$ estimators of the neutron spectrum. The fluctuations of the 100 spectra $`\stackrel{~}{x}`$ correspond to the propagation of the statistics fluctuations through the unfolding process.
$$\sigma _{j,statistics}^2=\frac{\underset{i=1}{\overset{N}{}}|\stackrel{~}{x}_{i,j}^2<\stackrel{~}{x}_{i,j}>^2|}{N}N=100$$
(37)
Also, to the 100 $`\stackrel{~}{x}`$ spectra correspond 100 reconstructed measurements $`\stackrel{~}{y}=A\stackrel{~}{x}`$. The reconstructed measurements $`\stackrel{~}{y}`$ are reinjected into the unfolding process giving new estimators of the spectra $`\stackrel{~}{x}_{new}`$ and twice reconstructed measurements $`\stackrel{~}{y}_{new}`$. The average of the absolute difference of the old and of the new spectra gives the systematics errors of the unfolding spectra :
$$\sigma _{j,systematics}^2=\frac{\underset{i=1}{\overset{N}{}}|\stackrel{~}{x}_{i,j}\stackrel{~}{x}_{i,j,new}|^2}{N}N=100$$
(38)
These two calculations of uncertainties are performed with 100 samples of Monte-Carlo generated data sets; this figure seems low due to the signal to noise ratio of 10 %, nevertheless the precision is sufficient for our purpose and moreover it is highly time consuming and we cannot afford greater computation. Also some tricks are used to avoid spending too much time recalculating 100 times almost the same things. Of course the $`\lambda `$ and $`\mu `$ parameters minimizations are very expensive on the CPU time balance, so the evolution with the iterations of the $`\lambda _n`$ (for example) is memorized during the first pass, while for the next passes the minimization of the $`\lambda `$ parameter is performed at the $`n^{th}`$ iteration in the neighboorhood of the $`\lambda _n`$ first evaluation. Another trick is used, oftenly after 50 or 80 iterations the $`\stackrel{~}{x}_n`$ and $`\stackrel{~}{y}_n`$ figures are stables due to the precision of the computation, so each successive iteration appears as an Identity operator, then at this moment the iterative process is stopped.
### 2.8 The initial guess
This kind of unfolding method shows several advantages. The most important one is the independance of the final solution regarding the initial guess. Thus any distribution can be used as an initial guess. We have tested several different initial guesses : random numbers, flat functions $`x_0=1`$ and the intercorrelation product between the response $`y`$ and the discrete response functions of matrix $`A`$.
$$x_{0,j}=<y_i,A_{i,j}\delta _j>\text{ }\delta _j\text{ : dirac function }$$
(39)
The strong point of this method is that using all these different initial guesses the iterative approach converges to exactly the same result. For a reason of fast and easy computation, we decided to use a flat initial guess : $`x_0=1`$.
### 2.9 Remarks about the minimization algorithms
To circonvene the problem of the minimization with $`\lambda `$ of the $`Q(\lambda )`$ criteria, we used different technics in order to maintain a high precision on the real $`\lambda `$ and a reasonable time of calculation for each iteration. The minimization of the criteria is performed in different ways to optimize this calculation time. Firstly, a multiple step scanning and secondly, a dichotomy algorithm are performed. The first solution of the first set $`y`$ is calculated carefully, while the other solutions computed to estimate the systematic and the statistic uncertainties, does begin their minimization at each step, with initial figures in the neighborhood of the first $`\lambda _n`$ sequence calculated previously. Of course, at this point the focusing on the right $`\lambda `$ is performed with scanning and dichotomy algorithms but with less steps, that avoid to spoil too much computation time.
## 3 Experimental approach
To illustrate this unfolding technics, we presents two examples of unfolding: the first one is a thermal neutron spectrum, the second one is a 565 keV mono-energetic neutron spectrum.
### 3.1 Unfolding a thermal neutron spectrum from the Sigma installation
This thermal neutron spectrum is produced near the Sigma facility at (IRSN-Cadarache), it consists of an assembling of several AmBe neutron sources surrounded by some graphite shielding. The expected spectrum is an heavy thermal component with some residual peaks between 1 and 3 MeV. Figures 1.a and 1.b shows precisely what was expected : a huge thermal neutron part and a small high energy component that goes from 100 keV up to 10 MeV. If the lower limit is physically explainable due to the tail of the decelarated neutrons from the high energy peaks, the upper bound around 10 MeV is not realistic and it is just due to the intrinsic resolution of the Bonner Sphere method itself. While figure 1.a shows the statistics error bars, figure 1.b presents the systematics unfolding error bars; in both case they are extremely tiny : the global statistics uncertainties is 0.11% while the global systematics errors of the unfolding is 3.6%. These very small uncertainties are mainly due to the very high number of counts registered for most of the sphere, this can be seen on figure 1.c, which represents the measured counts on the various sphere, (plain line), while the dot line corresponds to the recalculated spectrum after unfolding. These two shapes are in perfect agreement together. The last figures, shows one of the very specific parameter of the evolution of the unfolding process : $`1\lambda `$ vs. $`iteration`$ $`number`$. This shows that after only 20 iterations this parameter is already very low $`(10^{15})`$, that means that the iterative operator is very closed to the Identity operator, and the current solution is very closed to the final solution. This shape is very easy to unfold, this is why after only 78 iterations, the unfolding process is over. The main aspect of this unfolding is its ability to reveal some realistic details of the shape here a double structured shape with confortable uncertainties.
### 3.2 Unfolding a monoenergetic neutron spectrum at 565 keV
The Bonner Sphere system was operated in the beam of a Van de Graaf accelerator producing 565 keV mono-energetic neutrons at CEA (Bruyรจres le Chรขtel). The expected spectrum is a single energy peak enlarged by the energy resolution of the device eventually surrounded by some low energy diffused neutrons. The background of the low energy diffused neutrons is substracted from the measurements thanks to other measures undertook with copper-polyethylene shadow cones between the target that produces the neutrons and the detector. Figure 2c shows a light discrepancy between the measurements and the reconstructed data for the values between the 4 and the 7 inches spheres; this is probably due to the various efficiencies of the shadowing cones for small and big spheres. Anyway the global discrepancy between these two curves is 1.2 $`\%`$. Anywhere else, for the right and left wings, the reconstructed distribution is in good agreement with the original data. The unfolded spectrum (fig 2a and 2b), clearly shows a mono-energetic spectrum with a mean peaked at 686 keV. The difference between 565 keV and 686 keV can be partly explained by the large binning used for these energies because of the logarithmic progression of the bins in order to cover a wide energy spectrum of nine decades. Also the low energy resolution originated in the response matrix can be invoked to explain this discrepancy and finally too a systematic bias of about 40 keV from the energy calibration of this Van de Graaf generator. The shape of this spectrum is obviously peaked even if the FWHM is quite large : 726 keV. The global statistics errors represents 19. % of the unfolding spectrum (see figure 2a) while the systematics errors cover 28.8 % (see figure 2b); even if these two figures are quite large, this kind of spectrum estimation is sufficient enough for radioprotection purposes. These large values of uncertainties are probably due to the low counting rates compared to the thermal neutron measurement and also to the bad signal to background ratio for direct and diffused neutrons measures. Figure 2d illustrates the high speed of the unfolding process: after only 97 iterations, the stability is obtained.
### 3.3 Evolution of the minimization criteria $`Q_\lambda `$ for thermal neutron unfolding
The progress of the minimization criteria $`Q_\lambda `$ for the unfolding of the thermal neutron from Sigma is illustrated in figure 3. The $`Q_\lambda `$ vs. $`\lambda `$ function is drawn after 1, 2, 3 and finally 78 iterations. The evolution of this minimizing criteria is very fast. After the first iteration a wide minimum can be seen, while after the second iteration it is already very faint and extremly closed to 1. With the evolution of the iteration number the value of this minimum get smaller and smaller and also closer to 1. The value of $`Q_1`$ is always $`+\mathrm{}`$ this explains why the $`\lambda `$ minimum is very closed to 1 but never equals 1. So for the last iteration before stability, i=78, we find $`1\lambda 10^{16}`$ . This kind of unfolding scenario is typical of the cases we have studied, and also the order of magnitude of the last value of $`1\lambda 10^{16}`$.
## 4 Modification of the interpolation method for the thermal neutron part of the spectrum
### 4.1 Precise study of the thermal part of the unfolded spectrum of neutrons from the Sigma installation
If we consider into details the unfolded spectrum from the Sigma installation, we can be puzzled by the anomalous behaviour of this spectrum in the range $`10^810^7MeV(0.010.1eV)`$. Normally, the curve should look like a Maxwellian distribution with a temperature of 300K.
$`{\displaystyle \frac{dN}{dE}}E^{1/2}e^{\frac{E}{kT}}`$ $`kT=300K=0.025eV`$ (40)
Instead of this behaviour, the neutron spectrum distribution seems to rise very strightly at low energies without exhibiting a maximum around 25 meV. Three explanations can be argued for this fact. First of all, the low energy cut of the matrix response and of the unfolded spectrum is put at 10 meV $`(10^8eV)`$ which should distort a little bit this low energy region behaviour of thermal spectra. The main reason of this artificial cut at 10 meV is due to the very bad uncertainties of the response functions of the Bonner Sphere system to very low energy neutrons produced by the MCNP monte-carlo program in the range 1meV - 10 meV. So, very large statistical error bars are obtained through this method at 1 meV (bigger than 20%). This bad quality of results produced by MCNP for low energies neutrons is originated by the importance of the kinetic effect that polyethylen for example can generate through thermal exchanges between the molecules of polyethylene, whose agitation and binding energy are greater than 22 meV, and the impinging neutrons of very low energy 1 meV. Thus, these very cold neutrons are heated by the brownian motion of the molecules of the detector and a good deal of these neutrons do not penetrate inside the detector but are preferentially bouncing at the first contact with the surface of the sphere. This explanation is also at the origin of the discontinuities observed on the matrix response while the diameter of the spheres varies between the bare detector and the first 3 inches sphere for energies lower than 1 eV (see figure 4). The bouncing neutrons, due to thermal exchange between them and polyethylene explains the reheating and freezing neutrons between 1 meV and 1 eV. It produces for the spheres of diameter ranging between 1.3 to 2.5 inches at some energies an enhancement of the detectorโs efficiency and a reduction at other energies (see figure 5). This last point brings serious consequences while big discontinuities appears at low energies for small spheres between naked counter and the 3 inch sphere: we are not allowed to perform an interpolation of the measured data for hypothetical spheres in the range of the bare detector up to the 3 inch sphere. (Note : diameter of the naked counter =1.29 inch) This is equivalent to say that the measurement function is not smooth enough for spheres with diameters lower than 3 inches to authorize a direct interpolation. A solution would be to perform some more measurements with spheres of lower diameters than 3 inches, but due to these heavy discontinuities, small variations in diameter (less than 1mm) on the polyethylene thickness should generate strong variations on the neutronโs rates.
### 4.2 Rectifying the interpolation
We propose an alternative solution to this problem : in most of the actual cases, the low energy spectrum follows a Maxwellian distribution; the very rare cases where this assumption is wrong are very few and are well known situations as for example : a specific low energy neutrons filtered beam produced with a nuclear reactor adapted to solid state physics and cristallography. In other words, guessing a Maxwellian distribution is a low cost, extremly realistic and highly efficient assumption for low energies neutrons encoutered in actual radioprotection situations. So we performed an interpolation of the measurements for the spheres between 2.9 inches up to 12 inches, while we generated for a pure maxwellian distribution response of the spheres lower than 1.9 inches down to the naked counter, matching the data measured with this bare counter and its response to this pure Maxwellian distribution. For the hypothetical spheres between 1.9 and 2.9 inches, we performed a linear transition to connect these two curves of reconstructed measurements versus diameters of the spheres. Unfolding with this method the Sigma spectrum shows a satisfying solution regards to the Maxwellian distribution. The price to pay is a loss of continuity between the thermal and the epithermal parts of the unfolded spectrum. Anyway, the epithermal part is rather small compared to the pure thermal part and secondly the incidence of this epithermal part is rather small on the total dose flux measurement, because it deals with few neutrons ranging between 1 and 100 eV. The interesting point of this method is that it doesnโt affect at all the unfoldings of greater energies neutrons spectra, for example the 565 keV neutrons monoenergetic spectrum is absolutely unchanged!
### 4.3 Revisiting the Sigma thermal neutron unfolded spectrum in the light of the interpolation rectification
Using the rectification of the interpolation method with a Maxwellian distribution forcing for the thermal part of the spectra, onto a real data set from the Sigma thermal neutron facility, Am-Be sources + graphite moderator (see figure 6c). We obtain a very interesting neutron unfolded spectrum which exhibits a clear Maxwellian distribution for thermal neutrons of energy below $`10^7eV(100meV)`$, a high energy part of the spectrum for neutrons ranging from 1MeV up to 10 MeV corresponding to direct neutrons produced in the Am-Be source (see figure 6a). Between 10 and 100 eV $`(10^610^4MeV)`$, we can find an enhanced epithermal band, that seems to be disconnected from the Maxwellian thermal distribution. Here is the price to pay for this linear combination of two method of data interpolation (Spline+Maxwellian forcing) for the small diameters sphere countings.Anyway, this small epithermal band is not of a very importance for the dose flux calculation. Also, the global statistics errors remains very low: $`0.24\%`$ , compared to $`0.11\%`$ previously(see figure 6a). While the global systematics errors is slightly increased from $`3.6\%`$ to $`5.15\%`$ (see figure 6b). We can see that the final distribution is obtained after 85 iterations with a minimizing parameter $`1\lambda `$ less than $`10^{11}`$. The modification of the interpolation gives more realistic results when the original approach shows a small deficience due to discontinuities while interpolated. Fortunately we have a very good knowledge, thanks to the thermodynamic of the standard thermal distribution of low energies neutrons.
## 5 Conclusion
The development of this unfolding iterative Tykhonov approach was leaded by three important requirements : a high efficiency and fast converging method, a better and faster convergency than the least square method or than the classical Tykhonov method, but with the stability of the solution toward fluctuations of the measurements produced by the Tykhonov approach and also, a complete unsensitivity of the solution with respect to the initial guess. The iterated Tykhonov unfolding approach is modified to ensure a good convergence with a parameter running between 0 and 1. An accelerating algorithm is introduced between two classical Tykhonov iterations, this algorithm looks like a gradient method with a least square term. The calculation of statistics and systematics uncertainties is studied in details. The efficiency of this approach is presented in two experimental neutron spectrum unfoldings. This method is fully dedicated to Bonner Sphere system unfolding. Due to its very wide energy coverage, the Bonner sphere system is a very usefull device for radioprotection evaluation, but it cannot be used as a high resolution neutron spectrum instrument for metrology, even with an excellent unfolding method! One of the very big advantages of the method here illustrated is that no a priori information is introduced during the unfolding process leading to coherent results with the expected spectra. This independance toward the initial guess is extremely important for radioprotection measurement when a-priori no information on the neutron spectrum is available; not observing this point should lead to very prejudiciable wrong evaluations of dose flux calculation which is not allowed for radioprotection expertise approach. Nevertheless a very slight stretch to this law is applied in such a way that the interpolated counting rates between naked counter and the first 3 inches sphere is assumed to correspond to a pure Maxwellian distribution for the thermal neutrons part. This rectification is done to prevent from the high discontinuities of the matrix response for low diameters and low energies. This type of unfolding gives realistic results for thermal neutrons without modifying the high energies neutron spectra.
We must thank for all their help to build the Bonner Sphere, perform measurements, calibrations and matrix response calculation all the personal from the โLaboratoire de Recherche en Dosimรจtrie Externeโ of the IRSN and more specially Dominique Cutarella, Christian Itiรฉ and Emeric Pierre. We must also thank the PTB group of Germany and more specially M. Alevra who helped us to build this system C Bonner Sphere system. Most of this work was done at IRSN(Institut de Radioprotection et de Sรปretรฉ Nuclรฉaire (IRSN) Fontenay aux Roses (France) since 1997 up to 2002 at the Laboratory of External Dosemetry.
## 6 References
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# Nonclassical Moments and their Measurement
## I Introduction
Nonclassical effects have attracted substantial interest during the last decades. In particular due to the improvements of experimental techniques in the field of quantum optics nonclassical quantum states could be created in practice. After the first demonstration of photon antibunching Kimble et al. (1977) also sub-Poissonian photon statistics Short and Mandel (1983) and quadrature squeezing Slusher et al. (1985) could be realized.
In view of new possibilities of measurement and characterization of quantum states it became possible to more precisely consider the characterization of nonclassical effects. The experimental reconstruction of quantum states was demonstrated by balanced homodyne tomography Smithey et al. (1993) and some other methods, for a review see Welsch et al. (1999). In principle this allows one to completely characterize the quantum states of elementary quantum systems. On this basis an old question receives new interest: What are the typical signatures of nonclassical effects?
Besides balanced homodyne detection also homodyne correlation measurements have been considered. In the particular case of a weak local oscillator it has been shown that the method renders it possible to determine unusual types of moments, such as normally-ordered moments composed of both a field quadrature and the field intensity Vogel (1991, 1995); Carmichael et al. (2000). An important advantage of these measurement techniques consists in the fact that even for small quantum efficiencies the correlation properties of interest are not smoothed out, as it is known in balanced homodyning Welsch et al. (1999). For the following it is of some interest that the experimental determination of homodyne correlations in a multichannel device can be performed in practice Beck et al. (2001).
The definition of nonclassicality that is widely accepted in quantum optics is based on the existence of a well-behaved $`P`$-function. This means that a state is considered to have a classical counterpart if the $`P`$-function has the properties of a probability measure Titulaer and Glauber (1965); Mandel (1986), for a nonclassical state it fails to be interpreted as a probability. However, in addition to this definition also states have been considered to be nonclassical, whose mean photon number is small Mandel (1986) or whose quantum fluctuations are close to the vacuum-noise level Vogel (2000). It is also interesting that nonclassical effects in weak measurements have been considered to appear even for coherent states and for thermal states of small photon numbers Johansen (2004); Johansen and Luis (2004). The latter example would support both additional signatures of nonclassicality. However, the coherent-state example is only consistent with the requirement of the quantum noise being close to the vacuum level Vogel (2000), it does not require a small photon number. We also note that nonclassical states having negativities in the Wigner function are included in the class of nonclassical states we are dealing with.
Let us consider some recent developments of characterizing nonclassical effects of a single-mode quantum state. The attempt was made to formulate criteria that allow one to relate the nonclassicality to observable quantities. One approach consists in the use of characteristic functions of the quadrature distributions Vogel (2000). The usefulness of the nonclassicality criterion formulated in this way was successfully demonstrated in an experiment Lvovsky and Shapiro (2002). Later on this criterion turned out to be the lowest order of a hierarchy of necessary and sufficient conditions for nonclassicality Richter and Vogel (2002). This hierarchy is equivalent to the failure of Bochnerโs old criterion for the existence of a classical characteristic function Bochner (1933); Kawata (1972), when it is applied to the characteristic function of the $`P`$-function.
More recently the attempt was made to reformulate the notion of nonclassicality in a form which allows us to derive necessary and sufficient conditions in various representations Shchukin et al. (2005). For example, necessary and sufficient conditions in terms of quadrature moments have been obtained in this manner. These conditions in terms of moments are equivalent to those in terms of characteristic functions and to the failure of the $`P`$-function to be a probability distribution. The reformulation of the problem in terms of moments includes, as special cases, some previously discussed conditions Agarwal and Tara (1992); Agarwal (1993); Klyshko (1996). The connection of nonclassicality criteria with the $`17`$th Hilbert problem has also been considered Korbicz et al. (2005). Last but not least, it has been proposed to measure the nonclassicality of a single-mode quantum state via the entanglement potential Asbรณth et al. (2005). This is of particular interest since it shows that the further investigation of the nonclassicality of single-mode quantum states is of importance even for applications that require entangled state.
The aim of the present paper consists in a more general formulation of the conditions for the nonclassicality on the basis of observable moments Shchukin et al. (2005). Three different types of criteria will be analyzed in detail, which are formulated in terms of moments of annihilation and creation operators, of two quadrature operators, and of a quadrature and the number operator. In the latter case we will also discuss the relation of the criteria to Artinโs solution of the $`17`$th Hilbert problem. We will show that all the considered moments can be observed in a straightforward manner by homodyne correlation measurements. As a particular example for the usefulness of the methods under study we consider the characterization and detection of amplitude-squared squeezing Hillery (1987).
The paper is organized as follows. Necessary and sufficient criteria for the nonclassicality are derived in terms of appropriately chosen moments in Sec II. In Sec. III we propose some approaches for the measurement of the moments needed in the new versions of nonclassicality criteria. Special types of sufficient criteria for nonclassicality are consider in Sec. IV, with particular emphasis on the characterization of amplitude-squared squeezing. In Sec. V the concept is illustrated for a special type of minimum-uncertainty amplitude-squared squeezed states. A summary and some conclusions are given in Sec. VI.
## II Nonclassicality Criteria
The density operator $`\widehat{\varrho }`$ of any quantum state can be written in the diagonal coherent-state representation Sudarshan (1963); Glauber (1963)
$$\widehat{\varrho }=P(\alpha )|\alpha \alpha |d^2\alpha .$$
(1)
A quantum state is said to be nonclassical if the corresponding $`P`$-function (1) fails to be interpreted as a probability distribution on the complex plane Titulaer and Glauber (1965); Mandel (1986). Since the $`P`$-function may be highly singular, the criterion must be reformulated in terms of measured quantities before it could be applied for the interpretation of experiments.
In the following we will derive different versions of criteria for the nonclassicality of a quantum state in terms of normally-ordered moments of two appropriately chosen observables. It will be shown that among the possibilities of choosing two operators that completely describe the algebra of the harmonic oscillator there are at least two choices that lead to necessary and sufficient conditions for nonclassicality in terms of moments, for a brief consideration of one of these two cases see Shchukin et al. (2005). Other choices of operators, however, may only allow one to derive sufficient conditions for the nonclassicality in terms of moments. The reason for this is closely related to Artinโs solution of the $`17`$th Hilbert problem Prestel and Delzell (2001).
We will start with a brief review of the reformulation of the nonclassicality criteria in terms of characteristic functions Vogel (2000). In this manner it became possible to formulate a complete hierarchy of nonclassicality criteria that can be related to experimental data Richter and Vogel (2002). This approach will also be needed as the basis for a rigorous formulation of criteria in terms of moments.
### II.1 Characteristic Functions
The characteristic function $`\mathrm{\Phi }(\beta )`$ of $`P(\alpha )`$, that is its two dimensional Fourier transform, is defined as
$$\mathrm{\Phi }(\beta )=P(\alpha )e^{\alpha \beta ^{}\alpha ^{}\beta }d^2\alpha .$$
(2)
It is easy to verify that it obeys the conditions
$$\mathrm{\Phi }(0)=\mathrm{Tr}(\widehat{\varrho })=1,\mathrm{\Phi }(\beta )=\mathrm{\Phi }^{}(\beta ).$$
(3)
Moreover, $`\mathrm{\Phi }(\beta )`$ is a continuous function of $`\beta `$ for any quantum state. Therefore it obeys all the requirements to apply the following theorem introduced by Bochner Bochner (1933); Kawata (1972):
###### Theorem 1 (Bochner theorem)
: The function $`P(\alpha )`$ is a probability distribution on the complex plane if and only if for any smooth function $`f(\alpha )`$ with compact support the following expression is nonnegative
$$\mathrm{\Phi }(\alpha \beta )f^{}(\alpha )f(\beta )d^2\alpha d^2\beta 0.$$
(4)
The Bochner theorem can also be formulated in a discrete version by replacing the integrations in Eq. (4) with sums.
The necessary and sufficient conditions for $`P(\alpha )`$ being a probability measure can thus be reformulated in terms of the characteristic function. The relation (4) is fulfilled, if and only if for any order $`k=2,3,\mathrm{},`$ and for all complex $`\beta _1,\mathrm{},\beta _k`$ the determinants
$$D_k=D_k(\beta _1,\mathrm{},\beta _k)=\left|\begin{array}{cccc}1& \mathrm{\Phi }_{12}& \mathrm{}& \mathrm{\Phi }_{1k}\\ \mathrm{\Phi }_{12}^{}& 1& \mathrm{}& \mathrm{\Phi }_{2k}\\ 4\\ \mathrm{\Phi }_{1k}^{}& \mathrm{\Phi }_{2k}^{}& \mathrm{}& 1\end{array}\right|$$
(5)
obey the conditions
$$D_k0,$$
(6)
where $`\mathrm{\Phi }_{ij}=\mathrm{\Phi }(\beta _i\beta _j)`$. Consequently we arrive at the following theorem Richter and Vogel (2002):
###### Theorem 2
: A quantum state is nonclassical, if and only if there exist values $`\beta _i`$, ($`i=1,\mathrm{},k`$) for which at least one of the determinants $`D_k`$ ($`k=2,\mathrm{},\mathrm{}`$) attains negative values:
$$D_k<0.$$
(7)
We note that it is straightforward to relate the characteristic functions $`\mathrm{\Phi }(\beta )`$ to observable characteristic functions of quadrature distributions, for more details see Vogel (2000); Lvovsky and Shapiro (2002); Richter and Vogel (2002).
### II.2 Moments of $`\widehat{a}^{}`$, $`\widehat{a}`$
In order to derive criteria for nonclassicality in terms of moments, let us start with the normally-ordered expectation value of Hermitian operators of the form Shchukin et al. (2005); Korbicz et al. (2005)
$$:\widehat{f}^{}\widehat{f}:=|f(\alpha )|^2P(\alpha )d^2\alpha .$$
(8)
We will consider only such functions $`\widehat{f}=\widehat{f}(\widehat{a}^{},\widehat{a})`$ of the creation and annihilation operators, $`\widehat{a}^{}`$ and $`\widehat{a}`$, respectively, whose normally-ordered form exists. From the equation above it immediately follows that on a classical state the mean value (8) is nonnegative for any operator $`\widehat{f}`$. Hence the occurrence of negative mean values
$$:\widehat{f}^{}\widehat{f}:<0$$
(9)
is a clear signature of the nonclassicality of the quantum state under consideration.
Let us first derive conditions for classicality in terms of the moments $`\widehat{a}^k\widehat{a}^l`$. The fact that the mean value (8) is nonnegative for any polynomial function
$$\widehat{f}(\widehat{a}^{},\widehat{a})=\underset{k=0}{\overset{K}{}}\underset{l=0}{\overset{L}{}}c_{kl}\widehat{a}^k\widehat{a}^l$$
(10)
of $`\widehat{a}^{}`$ and $`\widehat{a}`$ leads to the nonnegativity of the following quadratic form
$$:\widehat{f}^{}\widehat{f}:=\underset{n,k=0}{\overset{K}{}}\underset{m,l=0}{\overset{L}{}}c_{kl}^{}c_{nm}\widehat{a}^{n+l}\widehat{a}^{m+k},$$
(11)
where the coefficients $`c_{nm}`$ are considered as independent variables. Due to Silvesterโs criterion it is equivalent to express this condition in terms of the determinants $`d_N`$ defined by
$$d_N=\underset{N}{\underset{}{\left|\begin{array}{ccccccc}1& \widehat{a}& \widehat{a}^{}& \widehat{a}^2& \widehat{a}^{}\widehat{a}& \widehat{a}^2& \mathrm{}\\ \widehat{a}^{}& \widehat{a}^{}\widehat{a}& \widehat{a}^2& \widehat{a}^{}\widehat{a}^2& \widehat{a}^2\widehat{a}& \widehat{a}^3& \mathrm{}\\ \widehat{a}& \widehat{a}^2& \widehat{a}^{}\widehat{a}& \widehat{a}^3& \widehat{a}^{}\widehat{a}^2& \widehat{a}^2\widehat{a}& \mathrm{}\\ \widehat{a}^2& \widehat{a}^2\widehat{a}& \widehat{a}^3& \widehat{a}^2\widehat{a}^2& \widehat{a}^3\widehat{a}& \widehat{a}^4& \mathrm{}\\ \widehat{a}^{}\widehat{a}& \widehat{a}^{}\widehat{a}^2& \widehat{a}^2\widehat{a}& \widehat{a}^{}\widehat{a}^3& \widehat{a}^2\widehat{a}^2& \widehat{a}^3\widehat{a}& \mathrm{}\\ \widehat{a}^2& \widehat{a}^3& \widehat{a}^{}\widehat{a}^2& \widehat{a}^4& \widehat{a}^{}\widehat{a}^3& \widehat{a}^2\widehat{a}^2& \mathrm{}\end{array}\right|}}.$$
(12)
The conditions for classicality than read as
$$d_N0.$$
(13)
To this end, however, these conditions have only been demonstrated to be necessary for the state to be classical. In order to show that they are necessary and sufficient, we will make use of Bochnerโs theorem.
To apply the Bochner theorem to moments, let us introduce the following operator $`\widehat{f}`$ for any smooth function $`f(\alpha )`$ with compact support:
$$\widehat{f}=\underset{ยฏ}{f}(\alpha ):\widehat{D}(\alpha ):d^2\alpha .$$
(14)
Due to the properties of $`f(\alpha )`$ and the following expansion of the normally-ordered displacement operator $`:\widehat{D}(\alpha ):`$,
$$:\widehat{D}(\alpha ):=\underset{k,l=0}{\overset{+\mathrm{}}{}}\frac{\alpha ^k(\alpha ^{})^l}{k!l!}\widehat{a}^k\widehat{a}^l,$$
(15)
the operator (14) is correctly defined and its normally ordered form
$$\widehat{f}=\underset{k,l=0}{\overset{+\mathrm{}}{}}c_{kl}\widehat{a}^k\widehat{a}^l$$
(16)
exists.
The left hand side of the expression (4) appearing in the Bochner theorem is nothing else but the mean value (8) for the operator (14),
$$:\widehat{f}^{}\widehat{f}:=\mathrm{\Phi }(\alpha \beta )\underset{ยฏ}{f}^{}(\alpha )\underset{ยฏ}{f}(\beta )d^2\alpha d^2\beta .$$
(17)
On the other hand the mean value $`:\widehat{f}^{}\widehat{f}:`$ can be represented as the following series
$$:\widehat{f}^{}\widehat{f}:=\underset{n,k,m,l=0}{\overset{+\mathrm{}}{}}c_{kl}^{}c_{nm}\widehat{a}^{n+l}\widehat{a}^{m+k}.$$
(18)
Suppose that all the determinants (12) are nonnegative. Then all finite sums (11) are also nonnegative. As the series (18) converges, it can be approximated by finite sums (11) and due to this it must be also nonnegative. Eventually the Bochner theorem states that this is equivalent to the nonnegativity of the $`P`$-function or the classicality of the state under consideration.
Consequently we formulate another theorem for the nonclassicality of a quantum state:
###### Theorem 3
: A quantum state is nonclassical if and only if at least one of the determinants $`d_N`$ violates the condition (13), that is
$$d_N<0,N=3,4,\mathrm{}.$$
(19)
Note that $`d_2`$ represent the incoherent part of the photon number,
$$d_2=\widehat{a}^{}\widehat{a}\widehat{a}^{}\widehat{a}.$$
(20)
Since this quantity is always nonnegative it yields no condition for nonclassicality.
It is possible to formulate other (sufficient) conditions for nonclassicality by considering subdeterminants of $`d_N`$. These subdeterminants are obtained by any pairwise cancellation of such lines and columns in $`d_N`$ that cross in a diagonal element of the matrix. The negativity of any such subdeterminant is a sufficient condition for nonclassicality. Criteria of this type may be useful for characterizing the nonclassical properties of special quantum states. Examples for such subdeterminants and for states that can be properly characterized by the resulting sufficient conditions will be studied in Sec. IV.
### II.3 Moments of $`\widehat{x}_\phi `$, $`\widehat{p}_\phi `$
Let us now consider two quadrature operators $`\widehat{x}_\phi `$ and $`\widehat{p}_\phi `$. As the quadrature operators are defined as linear combinations of $`\widehat{a}`$ and $`\widehat{a}^{}`$, the latter can be simply expressed in terms of the quadratures as
$$\widehat{a}=\frac{e^{i\phi }}{2}(\widehat{x}_\phi +i\widehat{p}_\phi ),\widehat{a}^{}=\frac{e^{i\phi }}{2}(\widehat{x}_\phi i\widehat{p}_\phi ).$$
(21)
One can reformulate the criteria for nonclassicality in terms of the normally-ordered moments $`:\widehat{x}_\phi ^n\widehat{p}_\phi ^m:`$ of the quadratures, as has been considered in Shchukin et al. (2005). We note that instead of using $`\widehat{x}_\phi `$ and $`\widehat{p}_\phi `$ the criteria could also be formulated with two arbitrary noncommuting quadratures $`\widehat{x}_\phi `$ and $`\widehat{x}_\phi ^{}`$, where $`\phi \phi ^{}\pm k\pi `$ ($`k=0,1,2,\mathrm{}`$). This more general form of the criteria is simply obtained by replacing $`\widehat{p}_\phi `$ with $`\widehat{x}_\phi ^{}`$ in the criteria derived below.
In the following we will make use of the fact that any operator $`\widehat{f}`$ whose normally-ordered form exists can be written as a normally-ordered power series with respect to $`\widehat{x}_\phi `$ and $`\widehat{p}_\phi `$,
$$\widehat{f}=\underset{n,m=0}{\overset{+\mathrm{}}{}}\stackrel{~}{c}_{nm}:\widehat{x}_\phi ^n\widehat{p}_\phi ^m:.$$
(22)
It is important that we use normally ordering here. An expansion of $`\widehat{f}`$ of the same structure, but the normally-ordered terms $`:\widehat{x}_\phi ^n\widehat{p}_\phi ^m:`$ being replaced with $`\widehat{x}_\phi ^n\widehat{p}_\phi ^m`$, may not exist. From the representation (22) it follows that the state is classical if and only if all the determinants $`\stackrel{~}{d}_N`$ defined by
$$\stackrel{~}{d}_N=\underset{N}{\underset{}{\left|\begin{array}{cccc}1& :\widehat{x}_\phi :& :\widehat{p}_\phi :& \mathrm{}\\ :\widehat{x}_\phi :& :\widehat{x}_\phi ^2:& :\widehat{x}_\phi \widehat{p}_\phi :& \mathrm{}\\ :\widehat{p}_\phi :& :\widehat{x}_\phi \widehat{p}_\phi :& :\widehat{p}_\phi ^2:& \mathrm{}\\ 4\end{array}\right|}}$$
(23)
are nonnegative:
$$\stackrel{~}{d}_N0.$$
(24)
Consequently, the criteria for nonclassicality can be equivalently formulated by the following theorem:
###### Theorem 4
: A quantum state is nonclassical if and only if at least one of the determinants $`\stackrel{~}{d}_N`$ violates the condition (24), that is
$$\stackrel{~}{d}_N<0,N=2,3,\mathrm{}.$$
(25)
In this version already the condition $`\stackrel{~}{d}_2<0`$ gives insight into the nonclassicality of quantum states. It represents the condition for quadrature squeezing,
$$\stackrel{~}{d}_2=:(\mathrm{\Delta }\widehat{x}_\phi )^2:<0.$$
(26)
Note that from the nonclassicality conditions in Eq. (25) together with (23) we may derive further conditions by pairwise cancellation of such lines and columns in the determinants that cross in the main diagonal. For example one may formulate conditions such as
$$s_\phi ^{(2)}=\left|\begin{array}{cc}:\widehat{x}_\phi ^2:& :\widehat{x}_\phi ^2\widehat{p}_\phi :\\ :\widehat{x}_\phi ^2\widehat{p}_\phi :& :\widehat{x}_\phi ^2\widehat{p}_\phi ^2:\end{array}\right|<0.$$
(27)
Conditions formulated in terms of such subdeterminants in some cases may be more efficient to describe particular nonclassical effects as the formal application of the hierarchy of the determinants $`d_N`$, for an example see Shchukin et al. (2005).
### II.4 Moments of $`\widehat{x}_\phi `$, $`\widehat{n}`$
It is also possible to use other pairs of operators that together with the unit operator generate the whole operator algebra of the harmonic oscillator. Let us consider the position operator $`\widehat{x}_\phi `$ and the photon number operator $`\widehat{n}`$. The annihilation and creation operators $`\widehat{a}`$ and $`\widehat{a}^{}`$ are expressed in terms of the $`\widehat{x}_\phi `$ and $`\widehat{n}`$ as follows
$$\widehat{a}=\frac{1}{2}(\widehat{x}_\phi +[\widehat{x}_\phi ,\widehat{n}])e^{i\phi },\widehat{a}^{}=\frac{1}{2}(\widehat{x}_\phi [\widehat{x}_\phi ,\widehat{n}])e^{i\phi }.$$
(28)
In the present case, however, we need to use the commutation relation $`[\widehat{a},\widehat{a}^{}]=1`$ to express the operators $`\widehat{a}`$, $`\widehat{a}^{}`$ in terms of $`\widehat{x}_\phi `$, $`\widehat{n}`$. Therefore it is no longer possible to substitute the expressions (28) into the expansion (16) and to rewrite the normally-ordered form (16) of the operator $`\widehat{f}`$ in terms of normally-ordered expressions in $`\widehat{x}_\phi `$ and $`\widehat{n}`$ as follows:
$$\widehat{f}=\underset{k,l=0}{\overset{+\mathrm{}}{}}\stackrel{~}{\stackrel{~}{c}}_{kl}:\widehat{x}_\phi ^k\widehat{n}^l:.$$
(29)
To make this more clear, already a linear term in (16), such as $`c_{01}\widehat{a}`$, has no representation in the form of Eq. (29).
As a consequence of this fact one cannot conclude that the determinants that are similar to (12) and (23), but with the moments $`:\widehat{x}_\phi ^k\widehat{n}^l:`$ instead of $`(\widehat{a}^{})^k\widehat{a}^l`$ and $`:\widehat{x}_\phi ^k\widehat{p}_\phi ^l:`$, respectively, form a complete hierarchy of criteria. The reason for this is the following. The not-existence of a representation of the form (29) for a general operator prevents one from relating these determinants to the Bochner theorem. Thus the corresponding determinants do not lead to necessary and sufficient conditions for nonclassicality.
Because of this situation we will formulate only necessary conditions for classicality or, the other way around, we will derive sufficient conditions for nonclassicality in terms of the moments $`:\widehat{x}_\phi ^k\widehat{n}^l:`$. The mean value of the normally-ordered operator $`:\widehat{F}(\widehat{x}_\phi ,\widehat{n}):`$ can be expressed in terms of the $`P`$-function as
$$:\widehat{F}(\widehat{x}_\phi ,\widehat{n}):=F(2\mathrm{Re}(\alpha e^{i\phi }),|\alpha |^2)P(\alpha )d^2\alpha .$$
(30)
For any classical state and for any operator $`F(\widehat{x}_\phi ,\widehat{n})`$ with a nonnegative associated c-number function $`F(2\mathrm{Re}\alpha ,|\alpha |^2)`$ the mean value (30) is nonnegative. If we identify $`:\widehat{F}(\widehat{x},\widehat{n}):`$ with $`:\widehat{f}^{}\widehat{f}:`$ together with the special form $`\widehat{f}`$ given in Eq. (29), we derive necessary conditions for classicality in terms of the determinants
$$d_N^{(1)}=\underset{N}{\underset{}{\left|\begin{array}{cccc}1& :\widehat{x}_\phi :& :\widehat{n}:& \mathrm{}\\ :\widehat{x}_\phi :& :\widehat{x}_\phi ^2:& :\widehat{x}_\phi \widehat{n}:& \mathrm{}\\ :\widehat{n}:& :\widehat{x}_\phi \widehat{n}:& :\widehat{n}^2:& \mathrm{}\\ 4\end{array}\right|}}$$
(31)
as
$$d_N^{(1)}0.$$
(32)
Consequently, if
$$d_N^{(1)}<0,N=2,3,\mathrm{},$$
(33)
the nonclassical properties of the considered quantum state have been verified by a sufficient but not necessary condition.
The latter can be illustrated as follows. There exist operators $`F(\widehat{x}_\phi ,\widehat{n})`$ of different kind whose associated $`c`$-number function $`\widehat{F}(2\mathrm{Re}(\alpha e^{i\phi }),|\alpha |^2)`$ is nonnegative. Let us consider
$$F(\widehat{x}_\phi ,\widehat{n})=(4\widehat{n}\widehat{x}_\phi ^2)\widehat{f}^{}(\widehat{x}_\phi ,\widehat{n})\widehat{f}(\widehat{x}_\phi ,\widehat{n}).$$
(34)
It is clear that the corresponding function is nonnegative
$$\begin{array}{cc}\hfill F(2\mathrm{Re}(\alpha e^{i\phi }),|\alpha |^2)& =4\mathrm{Im}^2\left(\alpha e^{i\phi }\right)\hfill \\ & \times |f(2\mathrm{Re}(\alpha e^{i\phi }),|\alpha |^2)|^20.\hfill \end{array}$$
(35)
This leads to the following necessary conditions for the nonnegativity of the $`P`$-function:
$$d_N^{(2)}0,$$
(36)
which is formulated in terms of the determinants
$$d_N^{(2)}=\underset{N}{\underset{}{\left|\begin{array}{ccc}4m_{01}m_{20}& 4m_{02}m_{21}& \mathrm{}\\ 4m_{02}m_{21}& 4m_{03}m_{22}& \mathrm{}\\ 3\end{array}\right|}}$$
(37)
with
$$m_{kl}=:\widehat{x}_\phi ^k\widehat{n}^l:.$$
(38)
Consequently, the violation of the nonnegativity of one such determinant,
$$d_N^{(2)}<0,N=1,2,\mathrm{},$$
(39)
is sufficient to demonstrate the nonclassical behavior of a quantum state. Note that these conditions already characterize nonclassical effects by the first order determinant $`d_1`$. By inspection of its expression,
$$d_1=4\widehat{n}:\widehat{x}_\phi ^2:=:\widehat{p}_\phi ^2:=:\widehat{x}_{\phi +\pi /2}^2:,$$
(40)
we observe that the condition $`d_1^{(2)}<0`$ reproduces the condition for quadrature squeezing provided that $`\widehat{x}_{\phi +\pi /2}=0`$, that is for quantum states whose mean amplitude is vanishing, as for example in a squeezed vacuum state.
### II.5 Relation to the 17th Hilbert problem
In his famous address to the 1900 International Congress of Mathematicians Hilbert formulated 23 problems that he considered to be the most important problems to be solved in the following century. The 17th problem, in its simplest form, was solved by Artin in 1926. Artinโs solution can be formulated in the following theorem (see e.g. Prestel and Delzell (2001)):
###### Theorem 5 ($`17`$th Hilbert problem)
: Any nonnegative polynomial $`F(x_1,\mathrm{},x_n)๐[x_1,\mathrm{},x_n]`$,
$$F(x_1,\mathrm{},x_n)0,(x_1,\mathrm{},x_n)๐^n,$$
(41)
can be represented as a finite sum of squares of rational functions
$$F(x_1,\mathrm{},x_n)=\underset{k=1}{\overset{N}{}}R_k^2(x_1,\mathrm{},x_n),$$
(42)
where $`R_k(x_1,\mathrm{},x_n)๐(x_1,\mathrm{},x_n)`$, i. e.
$$R_k(x_1,\mathrm{},x_n)=\frac{F_k(x_1,\mathrm{},x_n)}{G_k(x_1,\mathrm{},x_n)},$$
(43)
and all $`F_k`$ and $`G_k`$, $`k=1,\mathrm{},N`$ are polynomials
$$F_k(x_1,\mathrm{},x_n),G_k(x_1,\mathrm{},x_n)๐[x_1,\mathrm{},x_n].$$
(44)
This theorem is useful for getting deeper insight in the general problem of formulating conditions for the nonclassicality of quantum states in such cases as discussed in the preceding subsection. For this purpose let us consider the nonnegativity condition to be fulfilled for any quantum state that has a classical counterpart. Let us restrict the class of functions $`F(x_1,x_2)`$, with $`x_1=2\mathrm{Re}(\alpha e^{i\phi })`$ and $`x_2=|\alpha |^2`$, on the r.h.s. of Eq. (30) to polynomials. Now we may formulate the nonclassicality conditions in terms of the operators $`\widehat{x}_\phi `$ and $`\widehat{n}`$ by using Artinโs solution of the 17th Hilbert problem as follows:
$$:\widehat{F}(\widehat{x}_\phi ,\widehat{n}):=\underset{k=1}{\overset{N}{}}:\frac{\widehat{F}_k^2(\widehat{x}_\phi ,\widehat{n})}{\widehat{G}_k^2(\widehat{x}_\phi ,\widehat{n})}:<0.$$
(45)
That is, the condition for nonclassicality leads to normally-ordered expectation values of fractions of squares of polynomials. The extension of the conditions beyond polynomial functions is not obvious, nor it is straightforward to formulate necessary and sufficient conditions for nonclassicality in terms of normally-ordered moments of $`\widehat{x}_\phi `$ and $`\widehat{n}`$. Also other types of operators, even if they render it possible to span the algebra of the harmonic oscillator by using commutation rules, are expected to lead to similar problems.
In the cases of the operators $`\widehat{a}^{}`$ and $`\widehat{a}`$ on one hand and $`\widehat{x}_\phi `$ and $`\widehat{p}_\phi `$ on the other hand these problems do not exist. The general form of functions according to the solution of the 17. Hilbert problem is not needed due to the direct relation of the considered expectation values $`:\widehat{f}^{}\widehat{f}:`$ to the Bochner theorem. This allowed us to directly formulate necessary and sufficient conditions for the nonclassicality in terms of the moments of these operator. In the following we will consider the possibilities to measure these moments. This eventually allows one to characterize nonclassical states by measurable moments.
## III Measurement of moments
In the preceding section we have considered the formulation of criteria for the nonclassicality of quantum states in terms of moments. Such criteria are only of practical interest if one can design measurement principles that allow one to determine the corresponding moments. Since these schemes will depend on the types of moments, we will consider them separately. All these detection schemes will be based on homodyne correlation measurements. An important advantage of such measurements consists in the fact that the accessible information on a quantum state is not smoothed out be small detection efficiencies.
For the first time homodyne correlation measurements with a weak local oscillator have been proposed for measuring squeezing and anomalous moments, containing nonequal powers in the annihilation and creation operators, in resonance fluorescence Vogel (1991). This measurement principle was studied in more detail Vogel (1995), where the detection of quantum correlations of the photon number and a quadrature operator has been analyzed. Later on such measurements have also been considered for determining particular nonclassical properties of light Carmichael et al. (2000). In the following we will modify the method in such a way that it becomes possible to determine the different types of moments considered in the preceding section.
### III.1 Detection of $`\widehat{a}^k\widehat{a}^l`$
Let us consider the detection scheme presented in Fig. 1. As an example we have shown a measurement scheme of depth $`d=2`$, where the depth is the number of beam splitters between the entrance beam splitter $`\mathrm{BS}_0`$ and any of the detectors $`\mathrm{PD}_1,\mathrm{}\mathrm{PD}_4`$. An auxiliary photodetector $`\mathrm{PD}_\mathrm{a}`$ is used here and in the following measurement schemes to record, simultaneously with the correlation measurements of interest, the intensity of the local oscillator. In principle the measurement device could be further extended to an arbitrary value of depth $`d`$. Below we will show that such a scheme allows one to measure the moments $`\widehat{a}^k\widehat{a}^l`$ for $`k,l=0,1,\mathrm{},2^d`$.
The entrance beam-splitter for the signal field, $`\mathrm{BS}_0`$, is characterized by the parameters $`T_0`$ and $`R_0`$ with
$$|T_0||R_0|,$$
(46)
so that the signal field of interest is almost completely used for the correlation measurements. In our scheme the strength of the local oscillator on the detectors is typically of the same order of magnitude as the signal field, which is easily realized with a weakly reflecting entrance beam splitter. All the other beam splitters in the device are $`50`$%-$`50`$%, so that their parameters $`T_k`$ and $`R_k`$ can be written in the simple form
$$T_k=\frac{1}{\sqrt{2}}e^{i\phi _k},R_k=\frac{i}{\sqrt{2}}e^{i\phi _k}.$$
(47)
The operator $`\widehat{a}_0`$ describing the field behind the entrance beam splitter $`\mathrm{BS}_0`$ reads as
$$\widehat{a}_0=\frac{1}{\sqrt{2}}(T_0\widehat{a}+R_0\widehat{a}_{\mathrm{LO}}).$$
(48)
It is clear that all the operators $`\widehat{a}_k^{(d)}`$, that describe the fields on the $`k`$-th detector, are proportional to the operator $`\widehat{a}_0`$,
$$\widehat{a}_k^{(d)}=\frac{e^{i\mathrm{\Phi }_k}}{\sqrt{2^d}}\widehat{a}_0+\mathrm{},$$
(49)
where $`\mathrm{}`$ stands for a combination of the corresponding vacuum-channel operators which gives no contribution to the quantities considered below. The phase $`\mathrm{\Phi }_k`$ depends on the path leading from the beam-splitter $`\mathrm{BS}_0`$ to the $`k`$-th photodetector $`\mathrm{PD}_\mathrm{k}`$.
The measurement scheme can be used to detect the coincidences registered by all $`n`$ photodetectors, which are described by the normally-ordered correlation functions $`\mathrm{\Gamma }_{k_1,\mathrm{},k_n}`$ of the form
$$\mathrm{\Gamma }_{k_1,\mathrm{},k_n}=:\widehat{a}_{k_1}^{(d)}\widehat{a}_{k_1}^{(d)}\mathrm{}\widehat{a}_{k_n}^{(d)}\widehat{a}_{k_n}^{(d)}:.$$
(50)
Using Eq. (49) it is clear that these correlations do not depend on $`k_1,\mathrm{},k_n`$ and they can be written as
$$\mathrm{\Gamma }^{(n)}=\mathrm{\Gamma }_{k_1,\mathrm{},k_n}=\frac{1}{2^{nd}}\widehat{a}_0^n\widehat{a}_0^n.$$
(51)
Summing up over all possible combinations of $`n`$ photo-detectors, in order to avoid loss of measured data, we get the following function:
$$\begin{array}{cc}\hfill F_n& =\underset{\{k_1,\mathrm{},k_n\}}{}\mathrm{\Gamma }_{k_1,\mathrm{},k_n}=\hfill \\ & \left(\genfrac{}{}{0pt}{}{2^d}{n}\right)\mathrm{\Gamma }^{(n)}=\underset{m=n}{\overset{n}{}}f_n(m)e^{im\phi },\hfill \end{array}$$
(52)
where $`\phi =\phi _{\mathrm{LO}}\pi /2`$ and the coefficients $`f_n(m)`$ read
$$\begin{array}{cc}\hfill f_n(m)=\frac{\left(\genfrac{}{}{0pt}{}{2^d}{n}\right)}{2^{nd}}& \underset{kl=m}{}\left(\genfrac{}{}{0pt}{}{n}{k}\right)\left(\genfrac{}{}{0pt}{}{n}{l}\right)\times \hfill \\ & |T_0|^{k+l}|R_0\alpha |^{2nkl}\widehat{a}^k\widehat{a}^l.\hfill \end{array}$$
(53)
The function $`F_n`$ in Eq. (52) is a function of the phase $`\phi `$
$$F_n=F_n(\phi ),$$
(54)
and the coefficients $`f_n(m)`$ can be obtained from this function using Fourier transform
$$f_n(m)=\frac{1}{2\pi }_0^{2\pi }F_n(\phi )e^{im\phi }๐\phi .$$
(55)
Each coefficient $`f_n(m)`$ in Eq. (53) is a combination of some moments $`\widehat{a}^k\widehat{a}^l`$. From these combinations it is possible to extract the moments $`\widehat{a}^k\widehat{a}^l`$ themselves step by step. The moments $`\widehat{a}`$, $`\widehat{a}^{}`$ and $`\widehat{a}^{}\widehat{a}`$ can be obtained directly from $`F_1(\phi )`$,
$$\begin{array}{cc}\hfill \widehat{a}& =\frac{f_1(1)}{|T_0R_0\alpha |},\hfill \\ \hfill \widehat{a}^{}& =\frac{f_1(1)}{|T_0R_0\alpha |},\hfill \\ \hfill \widehat{a}^{}\widehat{a}& =\frac{f_1(0)|R_0\alpha |^2}{|T_0|^2}.\hfill \end{array}$$
(56)
Using these moments the next step gives the explicit expressions for the moments $`\widehat{a}^k\widehat{a}^l`$, $`k,l=1,2`$. We present the expressions only for the moments $`\widehat{a}^2`$ and $`\widehat{a}^{}\widehat{a}^2`$,
$$\begin{array}{cc}\hfill \widehat{a}^2& =\frac{8}{3}\frac{f_2(2)}{|T_0R_0\alpha |^2},\hfill \\ \hfill \widehat{a}^{}\widehat{a}^2& =\frac{1}{3}\frac{4f_2(1)3|R_0\alpha |^2f_1(1)}{|T_0^3R_0\alpha |}.\hfill \end{array}$$
(57)
Figure 2 illustrates the possibility to extract the general moments $`\widehat{a}^k\widehat{a}^l`$, $`0k,l2^d`$ step by step. The $`n`$-th step is the extraction of the moments $`\widehat{a}^k\widehat{a}^l`$ with $`0k,ln`$, i. e. the moments $`\widehat{a}^k\widehat{a}^l`$ with $`(k,l)`$ corresponding to the square
$$S_n=\{(k,l)๐^2|0k,ln\}.$$
(58)
But, in fact, we donโt measure the moments $`\widehat{a}^k\widehat{a}^l`$ themselves, we measure their combinations $`f_n(m)`$, $`m=n,\mathrm{},n`$, that consist of the moments $`\widehat{a}^k\widehat{a}^l`$, $`0k,ln`$ with additional condition $`kl=m`$. Geometrically the combination $`f_n(m)`$ contains the moments $`\widehat{a}^k\widehat{a}^l`$ with $`(k,l)`$ in the intersection of the square $`S_n`$ and the inclined line $`L_m:kl=m`$. This means that using the data obtained on one step only it is impossible to extract the moments $`\widehat{a}^k\widehat{a}^l`$, but assuming we have the data of previous steps it is: each combination $`f_{n+1}(m)`$ contains one new moment by comparison with $`f_n(m)`$. Note that $`f_{n+1}(n+1)`$ and $`f_{n+1}(n1)`$ contain one moment only, $`\widehat{a}^{n+1}`$ and $`\widehat{a}^{n+1}`$ correspondingly, so that these moments can be measured without keeping the information of previous steps. Geometrically this means that $`S_{n+1}L_m`$ contains one point more by comparison with $`S_nL_m`$. So, step by step it is possible to extract all the moments $`\widehat{a}^k\widehat{a}^l`$, $`0k,l2^d`$. The last step, the $`2^d`$-th, consists in the extraction of the moments $`\widehat{a}^k\widehat{a}^l`$ with $`k`$ or $`l`$ (or both) being equal to $`2^d`$.
### III.2 Detection of $`:\widehat{x}_\phi ^k\widehat{p}_\phi ^l:`$
The moments $`:\widehat{x}_\phi ^k\widehat{p}_\phi ^l:`$ can be measured with the scheme presented on the Fig. 3, as it has been proposed in Shchukin et al. (2005). In fact, this is a slightly modified version of the Noh-Fougรจres-Mandel device Noh et al. (1991). The lowest-order moments are explicitely expressed in terms of the correlation functions $`\mathrm{\Gamma }_j`$, $`\mathrm{\Gamma }_{jk}`$ according to the following formulas:
$$\widehat{x}_\phi =\sqrt{2}\frac{\mathrm{\Gamma }_3\mathrm{\Gamma }_4}{|\alpha |},\widehat{p}_\phi =\sqrt{2}\frac{\mathrm{\Gamma }_1\mathrm{\Gamma }_2}{|\alpha |}.$$
(59)
The second-order moments are of the form
$$\begin{array}{cc}\hfill :\widehat{x}_\phi ^2:& =4\frac{\mathrm{\Gamma }_{12}\mathrm{\Gamma }_{34}}{|\alpha |^2}+\widehat{n},\hfill \\ \hfill :\widehat{p}_\phi ^2:& =4\frac{\mathrm{\Gamma }_{34}\mathrm{\Gamma }_{12}}{|\alpha |^2}+\widehat{n},\hfill \\ \hfill :\widehat{x}_\phi \widehat{p}_\phi :& =4\frac{\mathrm{\Gamma }_{13}+\mathrm{\Gamma }_{24}\mathrm{\Gamma }_{12}\mathrm{\Gamma }_{34}}{|\alpha |^2}\widehat{n},\hfill \end{array}$$
(60)
where
$$\widehat{n}=\mathrm{\Gamma }_1+\mathrm{\Gamma }_2+\mathrm{\Gamma }_3+\mathrm{\Gamma }_4|\alpha |^2.$$
(61)
The amplitude $`|\alpha |`$ of the local oscillator is detected by the auxiliary photodetector $`\mathrm{PD}_\mathrm{a}`$. Note that one could also try to use the measured data more efficiently as discussed in detail in the preceding subsection. However, this would lead to somewhat more complex expressions for the moments. Since the procedure is straightforward, we do not present it here.
The scheme in the form shown in Fig. 3 allows one to determine all the moments $`:\widehat{x}_\phi ^k\widehat{p}_\phi ^l:`$ up to the order of $`k+l=4`$. The further extension of the method is straightforward. For example, each of the detectors $`\mathrm{PD}_i`$ ($`i=1\mathrm{}4`$) can be replaced with a beam splitter, each of which mixing the field with a vacuum input. In their output ports the output fields are detected by pairs of photodetectors. This allows us to extract the moments up to the order of $`k+l=8`$, and so forth.
### III.3 Detection of $`:\widehat{x}_\phi ^k\widehat{n}^l:`$
Figure 4 illustrates the possibility of measuring moments $`:\widehat{x}_\phi ^k\widehat{n}^l:`$ by an extended version of the homodyne cross correlation scheme considered in Ref. Vogel (1995). The moments $`\widehat{n}`$ and $`\widehat{x}_\phi `$ can be obtained according to the following relations
$$\begin{array}{cc}\hfill \widehat{n}& =\mathrm{\Gamma }_1+\mathrm{\Gamma }_2+\mathrm{\Gamma }_3+\mathrm{\Gamma }_4|\alpha |^2,\hfill \\ \hfill \widehat{x}_\phi & =\frac{1}{|\alpha |^2}\left(\mathrm{\Gamma }_1+\mathrm{\Gamma }_2\mathrm{\Gamma }_3\mathrm{\Gamma }_4\right).\hfill \end{array}$$
(62)
Higher-order moments can be extracted from the measured data step by step.
We give explicit expressions only for the second-order moments. The moment $`:\widehat{n}^2:`$ can be directly measured with blocked local oscillator. Note that this moment and the first-order moments together allow one to calculate the moment $`:\widehat{x}_\phi ^2:`$:
$$\begin{array}{cc}\hfill \mathrm{\Gamma }_{13}& +\mathrm{\Gamma }_{14}+\mathrm{\Gamma }_{23}+\mathrm{\Gamma }_{24}=\hfill \\ & \frac{1}{4}(:\widehat{n}^2:+2|\alpha |^2\widehat{n}+|\alpha |^4|\alpha |^2:\widehat{x}_\phi ^2:).\hfill \end{array}$$
(63)
The moment $`:\widehat{n}\widehat{x}_\phi :`$ can be obtained from the simple relation
$$\mathrm{\Gamma }_{12}\mathrm{\Gamma }_{34}=\frac{1}{4}(|\alpha |:\widehat{n}\widehat{x}_\phi :+|\alpha |^2\widehat{x}_\phi ).$$
(64)
All higher-order moments can be measured in such a way with an appropriately extended version of the scheme shown in Fig. 4.
## IV Amplitude-squared squeezing
In this section we will consider a particular nonclassical effect which can be well described in terms of the moments-based nonclassicality criteria derived in Sec. II. The examples will also be simple from the viewpoint of the measurement of the needed moments, when the measurement principles of Sec. III are used. This example is the amplitude squared squeezing Hillery (1987). We note that for the measurement of amplitude-squared squeezing to our knowledge there exists no direct measurement principle. It has been proposed to measure the effect after the interaction with a Kerr medium Hillery (1991). Such techniques are, however, not easy to use and require sufficiently strong signal fields for realizing the nonlinear interaction before the detection. Hence a possibility of a more direct measurement of the effect would be of interest.
Let us start with the characterization of the nonclassical effects we are interested in. As a special choice of $`\widehat{f}`$ in the nonclassicality condition (9) let us consider the following operator
$$\widehat{X}_\phi =\widehat{a}^2e^{i\phi }+\widehat{a}^2e^{i\phi }.$$
(65)
According to (8) for any number $`c`$ the mean value
$$:(\widehat{X}_\phi c)^2:0$$
(66)
is nonnegative for all classical states. For $`c=\widehat{X}_\phi `$ this gives the following condition:
$$:(\mathrm{\Delta }\widehat{X}_\phi )^2:0,$$
(67)
or, in other words, the condition
$$:(\mathrm{\Delta }\widehat{X}_\phi )^2:<0$$
(68)
is sufficient for nonclassicality. We note that the condition for amplitude-squared squeezing can be easily expressed as
$$s_\phi ^{(2)}=\left|\begin{array}{cc}1& :\widehat{x}_\phi \widehat{p}_\phi :\\ :\widehat{x}_\phi \widehat{p}_\phi :& :\widehat{x}_\phi ^2\widehat{p}_\phi ^2:& \end{array}\right|<0,$$
(69)
in terms of the moments of two quadratures.
Another way to formulate the condition for amplitude-squared squeezing is based on the moments of the operators $`\widehat{a}^{}`$, $`\widehat{a}`$. Let us consider the subdeterminant $`s_3`$ that results from the determinant (12) by canceling all its rows and columns except those beginning with the elements $`1`$, $`\widehat{a}^2`$ and $`\widehat{a}^2`$. This leads to a nonclassicality condition of the form
$$s_3=\left|\begin{array}{ccc}1& \widehat{a}^2& \widehat{a}^2\\ \widehat{a}^2& \widehat{a}^2\widehat{a}^2& \widehat{a}^4\\ \widehat{a}^2& \widehat{a}^4& \widehat{a}^2\widehat{a}^2\end{array}\right|<0.$$
(70)
We show that the negativity of this determinant is equivalent to the following condition
$$:(\mathrm{\Delta }\widehat{X}_\phi )^2:<0,\phi .$$
(71)
In fact, the determinant (70) can be written as follows:
$$s_3=\frac{1}{4}\underset{\phi }{\mathrm{min}}:(\mathrm{\Delta }\widehat{X}_\phi )^2:\underset{\phi }{\mathrm{max}}:(\mathrm{\Delta }\widehat{X}_\phi )^2:.$$
(72)
It is not difficult to see that the last term in the product on the right hand side of this equality is always nonnegative
$$\underset{\phi }{\mathrm{max}}:(\mathrm{\Delta }\widehat{X}_\phi )^2:0.$$
(73)
The explicit expressions for the minimum and the maximum of the variance of the operator $`\widehat{X}_\phi `$ are
$$\begin{array}{cc}\hfill \underset{\phi }{\mathrm{min}}:(\mathrm{\Delta }\widehat{X}_\phi )^2:& =2\left[\mathrm{\Delta }\widehat{a}^2\mathrm{\Delta }\widehat{a}^2\left|(\mathrm{\Delta }\widehat{a}^2)^2\right|\right],\hfill \\ \hfill \underset{\phi }{\mathrm{max}}:(\mathrm{\Delta }\widehat{X}_\phi )^2:& =2\left[\mathrm{\Delta }\widehat{a}^2\mathrm{\Delta }\widehat{a}^2+\left|(\mathrm{\Delta }\widehat{a}^2)^2\right|\right].\hfill \end{array}$$
(74)
It is clear that the mean value $`\mathrm{\Delta }\widehat{a}^2\mathrm{\Delta }\widehat{a}^2`$ is always nonnegative. Due to the second of the equalities (74) the maximal variance of $`\widehat{X}_\phi `$ is always nonnegative. Hence the condition for the negativity of $`s_3`$ in Eq. (72) is a direct demonstration of amplitude-squared squeezing.
An advantage of this method for demonstrating amplitude squared squeezing consists in the fact that we do not need to adjust the local oscillator phase to the noise minimum. The negativity of the considered subdeterminant demonstrates the negativity of the minimum of the normally-ordered variance already. Moreover, the moments $`\widehat{a}^2`$, $`\widehat{a}^4`$ and $`\widehat{a}^2\widehat{a}^2`$ needed in the condition (70) are easily obtained by the methods of the preceding section.
## V minimum-uncertainty amplitude-squared squeezed states
Let us consider amplitude squared-squeezed states in more detail (Hillery (1987); Bergou et al. (1991); Yu and Hillery (1994)). By definition, a state is amplitude-squared squeezed if
$$\underset{\phi }{\mathrm{min}}:(\mathrm{\Delta }\widehat{X}_\phi )^2:<0,$$
(75)
which can be rewritten in the equivalent form as
$$\phi :(\mathrm{\Delta }\widehat{X}_\phi )^2<4\widehat{n}+2.$$
(76)
We consider here a special class of amplitude-squared squeezed states that satisfy the minimum uncertainty relation Bergou et al. (1991). For this purpose we introduce the operators
$$\widehat{X}=\widehat{X}_\phi ,\widehat{Y}=\widehat{X}_{\phi +\pi /2}.$$
(77)
They fulfill the uncertainty relation
$$(\mathrm{\Delta }\widehat{X})^2^{1/2}(\mathrm{\Delta }\widehat{Y})^2^{1/2}4\widehat{n}+2.$$
(78)
A state satisfies the minimum uncertainty relation if
$$(\mathrm{\Delta }\widehat{X})^2^{1/2}(\mathrm{\Delta }\widehat{Y})^2^{1/2}=4\widehat{n}+2.$$
(79)
In the following we will only consider pure quantum states.
As it has been shown in Yu and Hillery (1994), a pure state $`|\psi `$ satisfies the condition (79) if it is a solution of the eigenvalue problem
$$(\widehat{X}+i\lambda \widehat{Y})|\psi =\beta |\psi ,$$
(80)
for all real nonnegative $`\lambda `$ and any complex $`\beta `$. One can easily check that form Eq. (80) it follows that
$$\begin{array}{cc}\hfill (\mathrm{\Delta }\widehat{X})^2& =\lambda (4\widehat{n}+2),\hfill \\ \hfill (\mathrm{\Delta }\widehat{Y})^2& =\frac{1}{\lambda }(4\widehat{n}+2).\hfill \end{array}$$
(81)
For either $`\lambda >1`$ or $`1/\lambda >1`$ one of the variations on the left-hand sides of Eqs (81) satisfies the condition (76). Therefore the state (80) is always amplitude-square squeezed, except for $`\lambda =1`$.
It was shown in Bergou et al. (1991) that there exist solutions of Eq. (80) of the form
$$|\psi (m,\lambda )=c_m(\lambda )\widehat{S}(z)H_m(i\gamma \widehat{a}^{})|0,$$
(82)
where
$$\gamma =\gamma (\lambda )=\{\begin{array}{cc}e^{i\pi /4}\sqrt{\sqrt{1\lambda ^2}/2}\hfill & 0<\lambda <1,\hfill \\ \sqrt{\sqrt{\lambda ^21}/2\lambda }\hfill & 1<\lambda ,\hfill \end{array}$$
(83)
and
$$|c_m(\lambda )|^2=\{\begin{array}{cc}1\hfill & m=0,\hfill \\ \frac{(1)^m}{m!C_m^m(2|\gamma (\lambda )|^2)}\hfill & m>0.\hfill \end{array}$$
(84)
The parameter $`\beta `$ in the equation (80) is connected with $`m`$ and $`\lambda `$ by the expression
$$\beta =\{\begin{array}{cc}i\sqrt{1\lambda ^2}(2m+1)\hfill & 0<\lambda <1,\hfill \\ \sqrt{\lambda ^21}(2m+1)\hfill & 1<\lambda .\hfill \end{array}$$
(85)
The parameter $`z=z(\lambda )=re^{i\phi }`$ reads
$$\mathrm{tanh}^2r=\frac{\lambda 1}{\lambda +1}e^{2i\phi },\phi =\{\begin{array}{cc}\pi /2\hfill & 0<\lambda <1,\hfill \\ 0\hfill & 1<\lambda .\hfill \end{array}$$
(86)
Some examples of the $`Q`$-function of the states (82) are shown in Fig. 5.
In Fig. 6 we illustrate examples of the determinants for the minimum-uncertainty amplitude-squared squeezed states under study. It is of great importance that all the moments appearing in the nonclassicality condition (70) can be determined by the homodyne correlation measurement technique proposed in Sec. III.1. Unless the methods of quantum-state reconstruction, for a review cf. Welsch et al. (1999), the measurements proposed here are even possible in cases when the overall quantum efficiency of the detection device is small.
## VI Summary and Conclusion
We have shown that condition for the nonclassicality of a single-mode quantum state, that is the failure of the $`P`$-function to be a probability measure, can be formulated in different ways. All the versions of necessary and sufficient conditions the we have formulated in this paper turn out to be special representations of the violation of Bochnerโs condition for the existence of a classical characteristic function of the $`P`$-function. The considered formulations of different forms of nonclassicality criteria include characteristic functions of quadratures and different types of normally-ordered moments. Most importantly, all the quantities under study are accessible to observation. Their measurement, however, requires to develop new types of measurement principles.
We have proposed rather simple and direct methods of measuring three kinds of normally-ordered moments of a single-mode radiation field. For each type of moments a particular detection scheme is analyzed: for the moments $`\widehat{a}^k\widehat{a}^l`$ of the creation and annihilation operators, the moments $`:\widehat{x}_\phi ^k\widehat{p}_\phi ^l:`$ of two quadratures and the moments $`:\widehat{x}_\phi ^k\widehat{n}^l:`$ of a quadrature and the photon number operator. In all the considered schemes the total number $`N`$ of photo-detectors needed to measure these moments is at most twice as large as the largest value of $`k`$ and $`l`$: $`N<2\mathrm{max}(k,l)`$. We also tried to present the extraction procedure for the moments in an optimal way form the viewpoint of an effective use of the available measured data. This means that one tries to use all the data obtained in measurements. Unfortunately this does not provide the expressions for the moments of interest in its simplest form. In the simplest form one would not need to sum up the correlation functions of all possible combinations of photo-detectors to get the moments. However, one would lose part of the measured data and thus this simplified approach would be less precise from the experimental point of view.
Our new methods of characterizing nonclassical effects by moments of annihilation and creation operators are applied to the characterization of amplitude-squared squeezing. We have shown that our criteria give a direct insight in the effect of amplitude-squared squeezing without the need of adjusting the phase of the local oscillator. More importantly, the demonstration of amplitude-squared squeezing until now was considered to require a rather difficult and indirect observation procedure, so that this effect received little attention in the context of experiments. The measurement procedures proposed in this paper turn out to be of particular simplicity for demonstrating the amplitude-squared squeezing effect. This may lead to new interest in this special nonclassical effect.
In conclusion we have formulated new types of criteria for characterizing nonclassical effects of radiation fields. For all versions of nonclassicality criteria we have proposed appropriate and simple measurement principles. These principles are based on homodyne correlation techniques with a weak local oscillator. They can be used even when the quantum efficiency of the device is small. The proposed methods may open new possibilities of demonstration and practical application of nonclassical effects of radiation fields.
###### Acknowledgements.
The authors gratefully acknowledge valuable discussions with R. Knรถrr, F. Liese and Th. Richter.
*
## Appendix A Moments of amplitude-squared squeezed states
The moments $`\widehat{a}^k\widehat{a}^l_m`$ of the state $`|\psi _m`$ can be obtained in the following way. Let us take
$$|\stackrel{~}{\psi }_m=c_m^1|\psi _m=\widehat{S}(z)H_m(i\gamma \widehat{a}^{})|0,$$
(87)
where the parameters $`z`$ and $`\gamma `$ are defined by the Eqs. (86) and (83) correspondingly. It is clear that
$$\widehat{a}^k\widehat{a}^l_m=\stackrel{~}{\psi }_m|\widehat{a}^k\widehat{a}^l|\stackrel{~}{\psi }_m.$$
(88)
We calculate the quantities $`\stackrel{~}{\psi }_m|\widehat{a}^k\widehat{a}^l|\stackrel{~}{\psi }_m`$ with the help of the generating function
$$F(x,y,u,v)=\underset{n,m,k,l=0}{\overset{+\mathrm{}}{}}\stackrel{~}{\psi }_n|\widehat{a}^k\widehat{a}^l|\stackrel{~}{\psi }_m\frac{x^ny^mu^kv^l}{n!m!k!l!}.$$
(89)
Clearly,
$$\begin{array}{cc}\hfill \underset{n,m=0}{\overset{+\mathrm{}}{}}& \stackrel{~}{\psi }_n|\widehat{a}^k\widehat{a}^l|\stackrel{~}{\psi }_m\frac{x^ny^m}{n!m!}=\hfill \\ & 0|e^{x^22i\overline{\gamma }x\widehat{a}}\widehat{S}^{}(z)\widehat{a}^k\widehat{a}^l\widehat{S}(z)e^{y^2+2i\gamma y\widehat{a}^{}}|0,\hfill \end{array}$$
(90)
and the function $`F(x,y,u,v)`$ reads as
$$\begin{array}{cc}& F(x,y,u,v)=\hfill \\ & 0|e^{x^22i\overline{\gamma }x\widehat{a}}\widehat{S}^{}(z)e^{u\widehat{a}^{}}e^{v\widehat{a}}\widehat{S}(z)e^{y^2+2i\gamma y\widehat{a}^{}}|0.\hfill \end{array}$$
(91)
Using the relation
$$\widehat{a}\widehat{S}(z)=\widehat{S}(z)(\mu \widehat{a}+\nu \widehat{a}^{})$$
(92)
it can be further simplified to the following form
$$F(x,y,u,v)=\mathrm{exp}\left(\frac{1}{2}๐ฐ^TA๐ฐ\right),$$
(93)
where $`๐ฐ=(x,y,u,v)^T`$ and
$$A=\left(\begin{array}{cccc}2& 4|\gamma |^2& 2i\overline{\nu }\gamma & 2i\mu \gamma \\ 4|\gamma |^2& 2& 2i\mu \overline{\gamma }& 2i\nu \overline{\gamma }\\ 2i\overline{\nu }\gamma & 2i\mu \overline{\gamma }& \mu \overline{\nu }& |\nu |^2\\ 2i\mu \gamma & 2i\nu \overline{\gamma }& |\nu |^2& \mu \nu \end{array}\right).$$
(94)
From this expression it follows that
$$\stackrel{~}{\psi }_n|\widehat{a}^k\widehat{a}^l|\stackrel{~}{\psi }_m=H_{nmkl}^{\{A\}}(0,0,0,0),$$
(95)
where $`H_{nmkl}^{\{A\}}(0,0,0,0)`$ is a four-dimensional Hermite polynomial defined as ($`๐=(a,b,c,d)^T`$)
$$\begin{array}{cc}\hfill \underset{n,m,k,l=0}{\overset{+\mathrm{}}{}}& H_{nmkl}^{\{A\}}(a,b,c,d)\frac{x^ny^mu^kv^l}{n!m!k!l!}=\hfill \\ & \mathrm{exp}\left(\frac{1}{2}๐ฐ^TA๐ฐ+๐ฐ^TA๐\right).\hfill \end{array}$$
(96)
Finally, we arrive at the result
$$\widehat{a}^k\widehat{a}^l_m=|c_m|^2H_{mmkl}^{\{A\}}(0,0,0,0),$$
(97)
for the moments we are interested in.
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# Spin relaxation in diluted magnetic semiconductor quantum dots
## I Introduction
The spin of the electron in low-dimensional semiconductor structures has received intense interest in recent years due to its potential applications in spintronic devices and quantum information processing technologies.Wolf ; Loss To improve the performance of such devices, the decoherence of the electron spin due to coupling to environmental degrees of freedom should be minimized. TheoreticalNazarov ; Nazarov2 ; Woods ; Efros ; Semenov and experimental investigationsAwschalom ; Paillard have shown that the electron spin could have a extremely long relaxation time in nonmagnetic semiconductor quantum dots (QDโs), compared with that in the bulk or quantum wells. Theoretical works proposed that the diluted magnetic semiconductor (DMS) QDโs can be used as spin aligners, spin memories as well as spin qubits.Loss2 DMS QDโs offer us a new flexibility in manipulating carrier spins, since the spin properties can be strongly influenced by applying an external magnetic field or varying the temperature.KChang Various relaxation mechanisms of the electron spin come from different coupling to the environment, i.e., magnetic impurity, nuclear spin, and spin-orbit interaction. It is important to identify the dominant mechanism of spin relaxation for a particular system. Previous theoreticalBastard ; Egues and experimentalAwschalom2 ; Akimoto works on spin relaxation in DMS quantum wells have indicated that the s-d exchange interaction between band electrons and localized spins of magnetic ions is the dominant spin-flip mechanism, leading to electron spin lifetime of the order of picoseconds. However, to the best of our knowledge, there is no study on the spin relaxation induced by the s-d exchange interaction in DMS QDโs. Detailed theoretical and experimental investigations are necessary to gain physical insight into the spin relaxation process in such structures.
In this paper, we investigate theoretically the spin relaxation of the lowest Zeeman doublet in vertical II-VI DMS QDโs. The spin-flip scattering caused by the acoustic phonon-mediated s-d exchange interaction between the conduction electron and Mn ions is considered. The electron-acoustic phonon coupling includes the piezoelectric and deformation potential interactions. Since the first-order spin-flip process through direct scattering by the Mn ions is generally blocked by the energy-matching condition, we consider the second-order process involving the emission or absorption of a phonon. Our calculation shows that this phonon-mediated spin-flip scattering leads to electron spin lifetime typically of the order of microseconds. The effectiveness of this mechanism increases significantly with increasing Mn concentration, electron spin splitting, vertical confining strength and lateral diameter, while it shows non-monotonic dependence on the magnetic field and temperature, due to the competing effect of the electron spin splitting, the phonon number and the correlation function of the Mn ions. It is interesting to notice that the spin relaxation of electrons in the lowest Zeeman doublet is suppressed in a strong magnetic field at low temperature for II-VI DMS QDโs with low Mn concentration.
The rest of this paper is organized as follows: the theoretical model and formula of the spin-flip scattering rate (SFR) induced by phonon-mediated s-d exchange interaction are derived in sec. II. Numerical results and discussions for the SFR as a function of magnetic field, as well as its dependence on the Mn concentration, QD size and temperature are given in sec. III and we give a brief conclusion in sec. IV.
## II Theory
We consider II-VI DMS QDโs subjected to a perpendicular magnetic field. Assuming an infinite deep well along the growth direction (the $`z`$ axis) and a in-plane parabolic confining potential, the electron wave function can be written as $`\psi (๐ซ)=\chi (z)\varphi (\rho ,\phi )`$, where $`\chi (z)=\sqrt{2/z_0}\mathrm{sin}(\pi z/z_0)`$ is the ground state wave function along the $`z`$ axis (we have assumed that the vertical confinement is strong and only the lowest energy level is relevant), $`z_0`$ is the width of the well, and $`\varphi (\rho ,\phi )`$ is the in-plane wave function determined by the two-dimensional Hamiltonian $`H=H_0+H_{sd}+H_{ep}`$. The first term
$$H_0=\frac{(๐ฉ+e๐)^2}{2m^{}}+\frac{1}{2}m^{}\omega _0^2\rho ^2+\frac{1}{2}g^{}\mu _BB\sigma _z$$
(1)
is the electron Hamiltonian in the external magnetic field and parabolic potential. Here $`m^{}`$ is the electron effective mass, $`๐=(By/2,Bx/2,0)`$ is the vector potential, $`\omega _0`$ characterizes the lateral confinement strength, $`g^{}`$ is the intrinsic electron g-factor, and $`\sigma _z`$ is the $`z`$-component of the Pauli matrices. The second term
$$H_{sd}=\underset{i}{}J(๐ซ๐_i)๐ฌ๐_i$$
(2)
describes the s-d exchange interaction between the electron ($`๐ฌ`$) and the localized Mn ion ($`๐_i`$), where $`J(๐ซ)`$ is the s-d coupling integral, and the summation runs over all the Mn sites. The last termNazarov2
$$H_{ep}=\underset{๐ช,\nu }{}\alpha _\nu (๐ช)(b_{๐ช,\nu }e^{i๐ช๐ซ}+b_{๐ช,\nu }^+e^{i๐ช๐ซ})$$
(3)
describes the interaction between the electron and acoustic phonon, where $`b_{๐ช\nu }(b_{๐ช\nu }^+)`$ is the annihilation (creation) operator of the bulk phonon mode with wave vector $`๐ช`$ and branch $`\nu `$.
The s-d exchange term is divided into a mean-field part and a fluctuating part, $`H_{sd}=H_{sd}^0+V_{sd},`$ where $`H_{sd}^0=\sigma _z\mathrm{\Delta }_{sd}/2,`$ $`V_{sd}=_iJ(๐ซ๐_i)(s^{(+)}S_i^{()}+s^{()}S_i^{(+)})/2,`$ $`s^{(\pm )}=s_x\pm is_y,`$ $`S_i^{(\pm )}=S_i^x\pm S_i^y,`$ $`\mathrm{\Delta }_{sd}=N_0\alpha xS_z`$ is the exchange splitting, $`N_0`$ is the number of unit cells per unit volume, $`\alpha =\varphi _c\left|J(๐ซ)\right|\varphi _c/\mathrm{\Omega }`$ ($`\mathrm{\Omega }`$ is the unit cell volume, $`\varphi _c`$ is the Bloch function at the bottom of the conduction band) is the s-d exchange coupling constant, $`x`$ is the fractional occupation factor of the cation sites by the Mn ions,
$$S_z=S_0B_S\left[\frac{g_{Mn}\mu _BBS}{k_B(T+T_0)}\right]$$
(4)
is the thermal average of the Mn spin, with $`S=5/2`$ the Mn 3d$`^\text{5}`$ spin, $`B_S(x)`$ the Brillouin function, and $`S_0,`$ $`T_0`$ phenomenological parameters accounting for the antiferromagnetic superexchange between neighboring Mn ions. Now the total Hamiltonian is divided into two parts, $`H=\overline{H}_0+V,`$ where
$$\overline{H}_0=\frac{p^2}{2m^{}}+\frac{1}{2}m^{}\omega ^2\rho ^2+\frac{1}{2}\omega _cL_z+\frac{1}{2}\sigma _z\mathrm{\Delta }_z,$$
(5)
$`V=V_{sd}+H_{ep}`$. Here $`\omega =\sqrt{\omega _0^2+\omega _c^2/4},`$ $`\omega _c=eB/m^{}`$ is the cyclotron frequency, $`L_z`$ is the $`z`$-component of the orbital angular momentum, $`\mathrm{\Delta }_z=g^{}\mu _BB+\mathrm{\Delta }_{sd}`$ is the total Zeeman splitting of the electron.
In order to obtain the SFR induced by the s-d exchange and electron-acoustic phonon interaction, we consider the whole system (including the electron, Mn ions and the phonon bath) transits from an initial state $`|i=|l\sigma ;M;N`$ to all possible final states $`|f=|l^{}\overline{\sigma };M^{};N^{}`$ in which the electron spin is reversed. Here $`|l\sigma `$ is the electron eigenstate ($`l`$ stands for the orbital quantum number $`(n,m)`$, see the Appendix for details), $`\sigma =\pm `$ ($`\overline{\sigma }=`$) denote spin-up (spin-down) and spin-down (spin-up) state, respectively. $`|M=|M_{1z},M_{2z},\mathrm{}`$ is the eigenstate of the Mn ions and $`|N=_{๐ช\nu }|n_{๐ช\nu }`$ denotes the phonon state. The transition rate is averaged over the random positions and the initial states of the Mn ions, as well as the initial states of the phonon system to give the SFR of the electron from the initial state $`|l\sigma `$ to the final state $`|l^{}\overline{\sigma }`$, denoted as $`W_{l^{}\overline{\sigma },l\sigma }`$.
Since the spin-flip process of electron is always accompanied by the flip of a Mn spin due to the conservation of the total angular momentum (see Eq. (2)), we introduce the renormalized electron energy $`E_{l\pm }=\epsilon _l\pm \mathrm{\Delta }_0/2`$, where $`\epsilon _l`$ is the orbital eigenenergy of $`\overline{H}_0`$ (see the Appendix for details) and $`\mathrm{\Delta }_0=\mathrm{\Delta }_{sd}\mathrm{\Delta }_i`$ is the (renormalized) electron spin splitting, with $`\mathrm{\Delta }_i=(g_{Mn}g^{})\mu _BB`$. Based on second-order perturbation theory, the transition amplitude between $`|i`$ and $`|f`$ is given by
$`T_{fi}`$ $`={\displaystyle \underset{l_1}{}}[{\displaystyle \frac{l^{}\overline{\sigma };M^{};N^{}|V_{sd}|l_1\sigma ;M;N^{}l_1\sigma ;M;N^{}|H_{ep}|l\sigma ;M;N}{E_{l^{}\overline{\sigma }}E_{l_1\sigma }}}`$ (6)
$`+{\displaystyle \frac{l^{}\overline{\sigma };M^{};N^{}|H_{ep}|l_1\overline{\sigma };M^{};Nl_1\overline{\sigma };M^{};N|V_{sd}|l\sigma ;M;N}{E_{l\sigma }E_{l^{}\overline{\sigma }}}}].`$
In the first term, the electron first hops from the initial state $`|l\sigma `$ to an virtual state with the same spin $`|l_1\sigma `$ through the interaction with a phonon, then it makes a spin-flip transition to the final state $`|l^{}\overline{\sigma }`$ through the s-d exchange interaction with one Mn ion. The second term describes the process that the electron is first scattered to an opposite-spin virtual state $`|l_1\overline{\sigma }`$ through the s-d exchange interaction with one Mn ion, then it transits to the final state $`|l^{}\overline{\sigma }`$ via the assistance of a phonon.
The scattering rate of the whole system is obtained from the Fermi golden rule $`W_{fi}=(2\pi /\mathrm{})\left|T_{fi}\right|^2\delta (E_fE_i)`$, and the SFR for the electron system is given by
$`W_{l^{},l+}`$ $`={\displaystyle \frac{1}{4}}x(N_0\alpha )^2G^+\left[n(\left|\mathrm{\Delta }_{ll^{}}\right|)+{\displaystyle \frac{1+\text{sign}(\mathrm{\Delta }_{ll^{}})}{2}}\right]K_{ll^{}},`$ (7)
$`W_{l+,l^{}}`$ $`={\displaystyle \frac{1}{4}}x(N_0\alpha )^2G^+\left[n(\left|\mathrm{\Delta }_{ll^{}}\right|)+{\displaystyle \frac{1\text{sign}(\mathrm{\Delta }_{ll^{}})}{2}}\right]K_{ll^{}},`$ (8)
where $`G^+=S^{()}S^{(+)},`$ $`G^+=S^{(+)}S^{()}`$ are correlation functions of the Mn ions, $`S^{(\pm )}=S_x\pm iS_y`$, $`\mathrm{\Delta }_{ll^{}}=\epsilon _l\epsilon _l^{}+\mathrm{\Delta }_0`$ is the electron energy detuning, $`n(E)=\left[\mathrm{exp}(E/(k_BT))1\right]^1`$ is the phonon distribution function, sign$`(x)=1`$ for $`x>0`$ and $`1`$ for $`x<0`$. The kernel $`K_{ll^{}}`$ is given by
$$K_{ll^{}}=\underset{l_1,l_2}{}\left[\frac{S_{l^{}l_1l^{}l_2}\mathrm{\Gamma }_{ll_1ll_2}}{\mathrm{\Delta }_{l_1l^{}}\mathrm{\Delta }_{l_2l^{}}}+\frac{S_{ll_1ll_2}\mathrm{\Gamma }_{l^{}l_1l^{}l_2}}{\mathrm{\Delta }_{ll_1}\mathrm{\Delta }_{ll_2}}\frac{2\mathrm{Re}(S_{l_1l_2ll^{}}\mathrm{\Gamma }_{l_1ll^{}l_2})}{\mathrm{\Delta }_{l_1l^{}}\mathrm{\Delta }_{ll_2}}\right],$$
(9)
where
$$S_{l_1l_2l_3l_4}=\mathrm{\Omega }d^3๐l_1|๐l_2|๐๐|l_3๐|l_4$$
(10)
is the dimensionless overlap integral, and
$$\mathrm{\Gamma }_{l_1l_2l_3l_4}=\frac{2\pi }{\mathrm{}}\underset{q,\nu }{}\left|\alpha _\nu (๐ช)\right|^2l_1|e^{i๐ช๐ซ}|l_2l_4|e^{i๐ช๐ซ}|l_3\delta (\mathrm{}\omega _{๐ช\nu }\left|\mathrm{\Delta }_{ll^{}}\right|)$$
(11)
is related to the spin-conserved phonon-induced transition rate. The explicit expressions for $`S_{l_1l_2l_3l_4}`$ and $`\mathrm{\Gamma }_{l_1l_2l_3l_4}`$ are given in the Appendix.
The spin lifetime $`\tau _{l\sigma }`$ of a given energy level $`|l\sigma `$ is given by
$$\frac{1}{\tau _{l\sigma }}=\underset{l^{}}{}W_{l^{}\overline{\sigma },l\sigma },$$
(12)
i.e., the sum of the spin-flip scattering rates from $`|l\sigma `$ to all opposite-spin final states $`|l^{}\overline{\sigma }`$.
## III Numerical results and discussions
From Eq. (9), we notice that the kernel $`K_{ll^{}}`$ and the SFR $`1/\tau _{l\sigma }`$ (see Eq. (12)) diverges when the energy of the intermediate state coincides with the initial or final state. To remove this divergence, we take into account the finite lifetime of the intermediate level and add a small broadening parameter (an order-of-magnitude estimate is 0.1meVBrunner ; Gammon ; Garcia ) to the energy of the intermediate state to convert this divergence into a resonance near the degeneracy point.Cardona This broadening parameter is not crucial for our calculation and would not change the qualitative behavior of the SFR.
We consider Cd<sub>1-x</sub>Mn<sub>x</sub>Te QDโs and use the following parameters in our numerical calculations: $`m^{}`$=0.096m<sub>0</sub> (m<sub>0</sub> is the free electron mass), $`g^{}`$=$``$1.6, CdTe mass density $`\rho `$=5.86 g/cm<sup>3</sup>, lattice constant $`a`$=0.6481 nm, $`g_{Mn}`$=2, $`S`$=5/2, $`N_0\alpha `$=220 meV, $`h_{14}`$=0.394$`\times `$10<sup>9</sup> V/m, sound velocity $`C_l`$=$`3083`$ m/s, $`C_t`$=1847 m/s. The lateral confining strength of the QD is characterized by the lateral diameter $`d=2\sqrt{\mathrm{}/(m^{}\omega _0)}`$. The dependence of $`S_0,`$ $`T_0`$ on the Mn concentration $`x`$ is obtained from Ref. 20.
Considering the electron occupies the lowest spin-up and spin-down levels $`(n=0,m=0,\pm )`$, i.e., the lowest Zeeman doublet in the DMS QD, the SFRโs can be calculated for the spin-up and spin-down states, which are denoted by $`1/\tau _+`$ and $`1/\tau _{}`$, respectively.
Since the spin-flip transitions to excited orbital levels are energetically unfavorable, the SFR of the lowest Zeeman doublet is dominated by the transition between the doublet, i.e., $`W_{00+,00}`$ and $`W_{00,00+}`$, such that the electron spin splitting $`\mathrm{\Delta }_0=\mathrm{\Delta }_{00,00}`$ and the kernel $`K_{00,00}`$ are important quantities for $`1/\tau _\pm `$ (see Eq. (7) and Eq. (8)). From Eq. (9), we note the contribution to $`K_{00,00}`$ comes mainly from the intermediate level whose energy is the closest to the initial (or final) state (i.e., the orbital state $`(0,1)`$ in most cases, see the Appendix for details). That is, the term containing $`S_{00,01,00,01}\mathrm{\Gamma }_{00,01,00,01}`$ and $`S_{01,01,00,00}\mathrm{\Gamma }_{00,01,01,00}`$ (from the Appendix, we see they are equal to each other, so both terms are denoted as $`S\mathrm{\Gamma }`$ for short) in Eq. (9) is the dominant contribution to $`K_{00,00}`$. Additional contributions to the SFR $`1/\tau _\pm `$ are the correlation function $`G^+`$, $`G^+`$, the phonon emission factor $`n(\left|\mathrm{\Delta }_0\right|)+1`$ or absorption factor $`n(\left|\mathrm{\Delta }_0\right|)`$. Therefore, in Fig. 1(a), (b), (c), (d), we plot schematically the correlation function $`G^+`$ and $`G^+`$, the phonon emission (absorption) factor and the product $`S\mathrm{\Gamma }`$ as a function of magnetic field, phonon energy, $`z_0`$and $`\left|\mathrm{\Delta }_0\right|,`$ respectively. In Fig. 1(a), we see that $`G^+=S(S+1)S_z^2\left|S_z\right|`$ decreases monotonically to zero while $`G^+=S(S+1)S_z^2+\left|S_z\right|`$ shows a peak and approaches a constant value with increasing magnetic field. Physically, this is because $`G^+`$ ($`G^+`$) is related to the transition of an electron from a spin-down (spin-up) initial state to a spin-up (spin-down) final state (see Eq. (7) and Eq. (8)). Due to the conservation of the total angular momentum, the $`z`$-component of the Mn spin $`S_z`$ should decrease (increase) by one in this process. In a strong magnetic field, however, all the Mn spins are polarized antiparallel the magnetic field (i.e., $`S_z5/2`$), thus the correlation function $`G^+`$ tends to vanish and the spin-flip process of the spin-down state is suppressed. The decrease of $`S\mathrm{\Gamma }`$ with increasing $`z_0`$ (approximately $`S\mathrm{\Gamma }1/z_0`$) and the peak behaviors of $`S\mathrm{\Gamma }`$ as a function of $`\left|\mathrm{\Delta }_0\right|`$ can be appreciated from Eq. (14), (15) and (16) in the Appendix. Note in the region where $`\left|\mathrm{\Delta }_0\right|`$ is small, the dependence of $`S\mathrm{\Gamma }`$ on $`\left|\mathrm{\Delta }_0\right|`$ is in agreement with Ref. 4. A peculiar feature is $`S\mathrm{\Gamma }`$ vanishes when the electron spin splitting $`\mathrm{\Delta }_0`$ approaches zero, due to the vanishing energy and, as a result, the vanishing density of states of the involved phonon.
Next, we shall investigate $`1/\tau _\pm `$ as a function of magnetic field for different temperatures, Mn concentrations and lateral diameters. The effect of vertical confining length $`z_0,`$ lateral diameter $`d`$ and temperature on $`1/\tau _\pm `$ is also presented.
### III.1 Strong lateral confinement
First we consider a small DMS QD with strong vertical confinement $`z_0`$=2 nm and small lateral diameter $`d`$=16 nm, such that the vertical and in-plane orbital energy separations are $``$3 eV and $``$12 meV, respectively. The large vertical orbital energy separation ensures only the lowest bound state is relevant to the spin relaxation, while the in-plane orbital energy separation which is much larger than $`\mathrm{\Delta }_i`$ ensures that the spin-flip transitions between the lowest Zeeman doublet usually dominate the spin relaxation process. However, if the Mn concentration is fairly high and the temperature is sufficiently low, the exchange splitting $`\mathrm{\Delta }_{sd}`$ may eventually become comparable with the orbital energy separation, then the lowest spin-up level may cross spin-down excited levels, opening up new spin relaxation channels for the spin-up state.
#### III.1.1 Low Mn concentration
The Mn concentration is taken as $`x`$=0.002, i.e., we take the saturated exchange splitting $`(\mathrm{\Delta }_{sd})_{sat}`$ $``$1 meV. The renormalized spin-dependent energy spectra for the electron at $`T`$=1 K and $`T`$=20 K are shown in Fig. 2. At low temperature $`T`$=1 K (Fig. 2(a)), the thermal-averaged Mn spin $`\left|S_z\right|`$ grows rapidly with increasing magnetic field (cf. Eq. (4)). As a result, the exchange splitting $`\mathrm{\Delta }_{sd}`$ increases rapidly to its maximum ($``$1 meV) and saturates, while $`\mathrm{\Delta }_iB`$ increases smoothly, such that the electron spin splitting $`\mathrm{\Delta }_0=\mathrm{\Delta }_{sd}\mathrm{\Delta }_i`$ reaches its maximum at a critical magnetic field, decreases when the magnetic field grows stronger, and eventually changes its sign at a strong enough magnetic field. In the high temperature case (see Fig. 2(b)), $`\left|S_z\right|`$ increases very slowly and $`\mathrm{\Delta }_i`$ always dominates, leading to a negative $`\mathrm{\Delta }_0`$ whose magnitude increases with increasing magnetic field or temperature.
The SFRโs of the lowest Zeeman doublet $`1/\tau _\pm `$ are shown in Fig. 3 as a function of magnetic field. The contributions from piezoelectric coupling (PZ) and deformation potential interaction (DP) are also indicated by the dashed and short-dashed lines, respectively. First, it is interesting to notice that $`1/\tau _{}`$ in Fig. 3(b) is significantly smaller than $`1/\tau _+`$ in Fig. 3(a), because the transition from $`(00)`$ to $`(00+)`$ needs to absorb a phonon, but the phonon number is very small at a low temperature $`T`$=1 K. Second, both $`1/\tau _+`$ and $`1/\tau _{}`$ are suppressed in a strong magnetic field at $`T`$=1 K, as shown in Fig. 3(a) and 3(b). The suppression of $`1/\tau _+`$ is due to the combined effect of small electron spin splitting $`\mathrm{\Delta }_0`$ (cf. Fig. 1(d)) and the vanishing phonon absorption factor $`n(\left|\mathrm{\Delta }_0\right|)`$, while the suppression of $`1/\tau _{}`$ is caused by the vanishing correlation function $`G^+`$ (see Fig. 1(a)). We also note that the PZ contribution dominates at small spin splitting $`(\mathrm{\Delta }_00.3`$ meV$`),`$ while the DP contribution dominates at large spin splitting $`(\mathrm{\Delta }_00.3`$ meV$`)`$, which can be clearly seen in Fig. 1(d), Fig. 3(c), 3(d) and all subsequent results. This is a direct result of the difference between the dependence of the PZ and DP coupling constant $`\alpha _\nu (๐ช)`$ on the wave vector: $`\alpha ^{PZ}(๐ช)1/\sqrt{q}`$, $`\alpha ^{DP}(๐ช)\sqrt{q},`$ such that the former (latter) dominates at small (large) phonon energy $`\mathrm{}\omega _๐ช(=\mathrm{\Delta }_0)`$. Finally, the electron spin splitting $`\mathrm{\Delta }_0`$ in Fig. 2(a) vanishes at $`B`$=0 T and $`B`$5 T. Correspondingly, the SFRโs $`1/\tau _\pm `$ vanish in both Fig. 3(a) and Fig. 3(b).
In the high-temperature regime, the spin splitting $`\left|\mathrm{\Delta }_0\right|`$ increases with increasing magnetic field. Consequently, the SFRโs $`1/\tau _\pm `$ exhibit the same behaviors, as shown in Fig. 3(c) and 3(d). Note, however, the contribution from PZ coupling decreases in a strong magnetic field. This is caused by the decreasing $`S\mathrm{\Gamma }`$ when the electron spin splitting $`\mathrm{\Delta }_0`$ exceeds $`\mathrm{\Delta }_{c1}`$ (see Fig. 1(d)). Compared with the low-temperature case, we find that the SFRโs increase significantly at a higher temperature, due to the increasing number of phonons. However, at very strong magnetic fields, the absence of high-energy phonon and the reduction of $`G^+`$, similar to the low-temperature case, reduce $`1/\tau _+`$ and $`1/\tau _{}`$, respectively.
#### III.1.2 Intermediate Mn concentration
The Mn concentration is increased to $`x`$=0.01, with a saturated exchange splitting $`(\mathrm{\Delta }_{sd})_{sat}`$ $``$5 meV, large enough to maintain a positive spin splitting $`\mathrm{\Delta }_0`$ in the whole range of the magnetic field $`B`$=0$``$8 T at low temperature. From the renormalized energy spectra shown in Fig. 4, we see that $`\mathrm{\Delta }_{sd}`$ ($`\mathrm{\Delta }_i`$) dominates at $`T`$=1 K ($`T`$=50 K) such that $`\mathrm{\Delta }_0`$ is positive (negative) over the whole range of the magnetic field. In Fig. 5, the SFRโs $`1/\tau _\pm `$ are plotted as a function of magnetic field for $`T`$=1 K and $`T`$=50 K, respectively. At $`T`$=1 K, when the magnetic field increases, $`1/\tau _+`$ increases to its saturation value due to the increase and saturation of $`\mathrm{\Delta }_{sd}`$ (and thus $`\mathrm{\Delta }_0`$), while that of the spin-down level is much smaller and shows a sharp peak at very weak magnetic field $`(B`$0.04 T$`)`$ and reduces to zero quickly, which is caused primarily by the small phonon absorption factor $`n(\left|\mathrm{\Delta }_0\right|)`$ at low temperature and partly by the reduction of the correlation function $`G^+`$ in a strong magnetic field.
The further increase of $`1/\tau _+`$ with increasing magnetic field in Fig. 5(a) is due to the decrease of the orbital excitation energy (in Eq. (9), the denominator $`\mathrm{\Delta }_{l_1l_2}`$ consists of two parts, the orbital excitation energy $`\epsilon _{l_1}\epsilon _{l_2}`$ and the spin splitting $`\mathrm{\Delta }_0`$). At $`T`$=50 K, $`1/\tau _+`$ and $`1/\tau _{}`$ both increase with increasing magnetic field, due to the increase of the spin splitting $`\mathrm{\Delta }_0`$and the number of phonons. Note in Fig. 5(c) and 5(d), the zero-field SFRโs do not vanish. This can be understood because the electron can transit to excited orbital levels at $`T`$=50 K, although the spin-flip transition to the ground orbital level is prohibited due to vanishing spin splitting $`\mathrm{\Delta }_0`$. Furthermore, we notice the SFR in the case of intermediate Mn concentration $`(x`$=0.01$`)`$ is several times larger than that for low Mn concentration $`(x`$=0.002$`),`$ since $`1/\tau _\pm `$ $`x`$ through Eq. (7) and Eq. (8) (note, however, the effect of the Mn concentration $`x`$ is also manifested through changing the electron spin splitting $`\mathrm{\Delta }_0`$). Finally, the decrease of the PZ contribution at large electron spin splitting $`\mathrm{\Delta }_0`$, and the resulted crossing of the PZ and DP contributions at $`\mathrm{\Delta }_00.3`$ meV, as discussed in the previous subsection, is again observed.
#### III.1.3 High Mn concentration
In this subsection, the Mn concentration is increased further to $`x`$=0.05, with a saturated exchange splitting $`(\mathrm{\Delta }_{sd})_{sat}`$ $``$16 meV, such that $`\mathrm{\Delta }_{sd}`$ is comparable with the in-plane orbital level separation ($`12`$ meV). The renormalized energy spectra at $`T`$=1 K and $`T`$=50 K are shown in Fig. 6. We see from the left panel ($`T`$=1 K) that the spin-up ground orbital level crosses the spin-down excited orbital level at a critical magnetic field $`B_c`$2.2 T. At a higher temperature $`T`$=50 K, $`\mathrm{\Delta }_{sd}`$ is still the dominant contribution to $`\mathrm{\Delta }_0,`$ but its magnitude decreases, such that the energy levels do not cross (see Fig. 6(b)).
The SFRโs of the lowest Zeeman doublet are shown in Fig. 7. First we note the SFRโs in panel (a), (c), (d) are of the same order of inverse nanoseconds, while that in panel (b) is much smaller, due to the absence of high-energy phonons at low temperature. The most significant feature is the sharp peak around the critical magnetic field $`B_c`$ in Fig. 7 (a), corresponding to the crossing of the spin-up ground level with the first excited spin-down level (see Fig. 6(a)). This is because the level crossing leads to a resonance in $`K_{00,00}`$ and, as a result, in $`1/\tau _+`$ (see the discussion in the beginning of section III).
Note, however, $`1/\tau _{}`$ in Fig. 7(b) doesnโt show this resonant behavior, because the resonance of $`K_{00,00}`$ is suppressed by the vanishing phonon absorption factor $`n(\left|\mathrm{\Delta }_0\right|)`$. For $`B<B_c,`$ the transition to the lowest spin-down level gives the dominant contribution to $`1/\tau _+`$, which reaches its maximum at $`B`$0.8 T and decreases at stronger magnetic fields (cf. Fig. 1(d)). For $`B>B_c,`$ the electron in the lowest spin-up level can transit into the first excited spin-down level, opening up a second spin-flip channel, and it is just the contribution from this channel that dominates in the $`B>B_c`$ regime. This second contribution reaches its maximum value at $`B`$4 T and then decreases, which can also be interpreted via Fig. 1(d). For the $`T`$=50 K case, $`1/\tau _+`$ and $`1/\tau _{}`$ both increase with increasing magnetic field, showing a broad peak at $`B`$2.5 T. The peak comes from the competing effect of increasing spin splitting $`\mathrm{\Delta }_0`$ (which leads to increasing $`1/\tau _\pm `$) against decreasing correlation function $`G^+,G^+`$ and phonon emission (absorption) factor.
### III.2 Weak lateral confinement
Now we turn to investigate QDโs with weak lateral confinement $`d`$=40 nm, whose orbital level separation is comparable with $`\mathrm{\Delta }_i`$. In this case, with small Mn concentration or high temperature, $`\mathrm{\Delta }_i`$ makes the main contribution to the spin splitting. Consequently, the spin splitting $`\mathrm{\Delta }_0`$ is negative, and the lowest spin-down level may cross the excited spin-up levels in a strong magnetic field. Figure 8 shows the renormalized electron energy spectra at $`T`$=1 K and $`T`$=10 K.
In panel (a), where the temperature is low, the spin-down ground level first crosses the spin-up ground level, due to the small exchange splitting $`\mathrm{\Delta }_{sd}`$ and low temperature, then it sweeps cross the excited spin-up levels, due to the large $`\mathrm{\Delta }_i`$ compared with the orbital excitation energy in a strong magnetic field. When the temperature increases to $`T`$=10 K (see Fig. 8(b)), the exchange splitting is suppressed and $`\mathrm{\Delta }_i`$ always dominates. The lowest spin-down level crosses the excited spin-up levels but the crossing between the Zeeman split doublet doesnโt occur.
Figure 9 shows $`1/\tau _\pm `$ as a function of magnetic field at $`T`$=1 K and $`T`$=10 K, respectively.
In panel (a) and (b), the low-field behaviors of $`1/\tau _\pm `$ resemble those of strong lateral confinement and low Mn concentration (see Fig. 3(a) and 3(b)). At higher temperature, $`1/\tau _\pm `$ exhibit many peaks at higher magnetic fields (indicated by the arrows), which are caused by the aforementioned level crossings. However, in very strong magnetic fields, the peaks are suppressed by the phonon absorption factor (for $`1/\tau _+`$) and the correlation function $`G^+`$ (for $`1/\tau _{}`$). At a higher temperature $`T`$=10 K, the resonances of the kernel $`K_{00,00}`$ are less suppressed and more pronounced peaks arises in $`1/\tau _\pm `$, leading to short spin lifetimes of the order of nanoseconds, compared with the microsecond scale in the $`T`$=1 K case.
### III.3 Temperature effect
In the above, we have observed that the temperature plays an important role in determining the SFR through changing the electron spin splitting $`\mathrm{\Delta }_0`$, the correlation function $`G^+`$, $`G^+`$ and the phonon emission (absorption) factor. Taking a small QD ($`z_0`$=2 nm, $`d`$=16 nm) for example, we plot the renormalized energy spectrum and $`1/\tau _\pm `$ as a function of temperature at $`B`$=4 T in Fig. 10.
It is interesting to notice that the Zeeman split doublet crosses each other at an elevated temperature (see Fig. 10(a)), due to the reduction of the exchange splitting $`\mathrm{\Delta }_{sd}`$. From Fig. 1(d), we see that the quantity $`S\mathrm{\Gamma }`$ decreases with decreasing $`\mathrm{\Delta }_0`$ for small $`\mathrm{\Delta }_0`$. This effect, together with the phonon emission factor, which shows a sharp peak at $`\mathrm{\Delta }_0`$=0, results in the non-monotonic temperature dependence of $`1/\tau _+`$ shown in Fig. 10(b). In Fig. 10(c), the low-temperature SFR for the spin-down level $`1/\tau _{}`$ vanishes due to the absence of high-energy phonons and the vanishing correlation function $`G^+`$. Note at $`T`$34 K, the SFR of both spin-up and spin-down levels vanishes, due to the vanishing spin splitting $`\mathrm{\Delta }_0`$.
### III.4 Dependence of the SFR on the confinement
Both the vertical and lateral confinement of the QD can affect the spin relaxation significantly through varying the electron wave function. The effect of vertical confinement on $`1/\tau _\pm `$ comes from the form factor $`Z(q)`$ (see Eq. (17)) and the overlap integral (see Eq. (14)). It can be seen from Fig. 1(c) that the quantity $`S\mathrm{\Gamma }`$ is roughly proportional to $`1/z_0,`$ such that the SFRโs should also show the same behavior, which can be seen in Fig. 11(a) and 11(b). Note here the largest $`z_0`$ (20nm) still sustains a vertical orbital level separation of $`30`$ meV, such that the influence of higher subbands on the spin relaxation is negligibly small. The approximate relationship $`1/\tau _\pm 1/z_0`$ comes from two factors. First, the s-d exchange scattering amplitude with one Mn ion is proportional to $`1/z_0`$ which, when squared (in the Fermi golden rule) and averaged over all the Mn sites, leads to the $`1/z_0`$ dependence of the spin-flip scattering rate to a given final state. Second, the spin-flip channel (i.e., the number of final states) doesnโt increase provided $`z_0`$ is small enough such that only the lowest bound state is relevant. We notice that G. Bastard et al. performed a theoretical calculation of the SFR of subbands in DMS quantum wells, and similar dependence of the SFR on the well width is predicted.Bastard
The effect of the lateral confinement strength, characterized by the lateral diameter $`d`$, on the spin relaxation is shown in Fig. 11(c) and 11(d). The $`d^4`$ dependence of $`1/\tau _\pm `$ can be appreciated as follows. From the Appendix and the definition $`d=2\sqrt{\mathrm{}/(m^{}\omega _0)}`$, we see the orbital level separation $`\delta `$ is roughly proportional to $`1/d^2`$, while $`1/\tau _\pm `$ is inversely proportional to $`\delta ^2`$ (see Eq. (9)) when the spin splitting is small, so we expect that $`1/\tau _\pm `$ should be approximately proportional to $`d^4`$, although the precise dependence of $`1/\tau _\pm `$ on $`d`$ is also affected by the phonon-induced transition rate $`\mathrm{\Gamma }`$ (cf. Eq. (15), (16)).
The dependence of the SFR on the QD size $`1/\tau _\pm z_0^1d^4`$ caused by phonon-mediated s-d exchange scattering is quite different from those caused by other spin relaxation mechanisms in nonmagnetic semiconductor QDโs.Woods ; Smimov ; MWWu We note that this relationship can be deduced from the work by Nazarov,Nazarov2 where the spin relaxation in nonmagnetic semiconductor QDโs is considered, but the magnitude of the SFR in our results is several orders of magnitude higher.
## IV Conclusions
Based on second-order perturbation theory, we have investigated the SFR caused by the phonon-mediated s-d exchange interaction of the lowest Zeeman split doublet in II-VI DMS QDโs as a function of magnetic field, as well as the dependence of the SFR on the Mn concentration, dot size and temperature. We found the SFR increases with increasing Mn concentration and electron spin splitting $`\mathrm{\Delta }_0`$. Increasing the lateral dot size leads to enhanced SFR while increasing the vertical dot size reduces the SFR for a small QD. The dependence of the SFR on the magnetic field and temperature shows non-monotonic behaviors, due to the competing effect between the electron spin splitting, the phonon emission (absorption) factor and the correlation function of the Mn ions. It is interesting to notice that the spin relaxation of both spin-up and spin-down electrons is suppressed in the case of strong magnetic field and low Mn concentration at low temperature.
###### Acknowledgements.
This work was supported by the NSFC No. 60376016, 863 project No. 2002AA31, and the special fund for Major State Basic Research Project No. G001CB3095 of China.
*
## Appendix A EXPRESSIONS OF $`S`$ AND $`\mathrm{\Gamma }`$
The orbital part of $`\overline{H}_0`$ gives the Fock-Darwin states
$`\varphi _{nm}(\rho ,\phi )`$ $`={\displaystyle \frac{1}{\sqrt{2\pi }}}e^{im\phi }R_{nm}(\rho )(n,|m|=0,1,2,\mathrm{}),`$ (13)
$`R_{nm}(\rho )`$ $`={\displaystyle \frac{\sqrt{2}}{l_0}}\sqrt{{\displaystyle \frac{n!}{(n+\left|m\right|)!}}}({\displaystyle \frac{\rho }{l_0}})^{\left|m\right|}\mathrm{exp}({\displaystyle \frac{\rho ^2}{2l_0^2}})L_n^{\left|m\right|}({\displaystyle \frac{\rho ^2}{l_0^2}}),`$
with corresponding orbital energy $`\epsilon _{nm}=(2n+\left|m\right|+1)\mathrm{}\omega +m\mathrm{}\omega _c/2,`$ where $`l_0=\sqrt{\mathrm{}/(m^{}\omega )},`$ $`L_n^m(x)`$is the generalized Laguerre polynomial. For convenience, we also introduce $`n_+=n+(\left|m\right|+m)/2,`$ $`n_{}=n+(\left|m\right|m)/2.`$
The dimensionless overlap integral is
$`S_{l_1l_2l_3l_4}`$ $`=\delta _{m_1+m_2,m_3+m_4}\sqrt{{\displaystyle \frac{n_1!n_2!n_3!n_4!}{(n_1+\left|m_1\right|)!(n_2+\left|m_2\right|)!(n_3+\left|m_3\right|)!(n_4+\left|m_4\right|)!}}}{\displaystyle \frac{\mathrm{\Omega }\xi }{\pi z_0l_0^2}}`$ (14)
$`{\displaystyle _0^{\mathrm{}}}e^{2x}(\sqrt{x})^{\left|m_1\right|+\left|m_2\right|+\left|m_3\right|+\left|m_4\right|}L_{n_1}^{\left|m_1\right|}(x)L_{n_2}^{\left|m_2\right|}(x)L_{n_3}^{\left|m_3\right|}(x)L_{n_4}^{\left|m_4\right|}(x)๐x,`$
where $`\xi =z_0๐z\left|\chi (z)\right|^4=3/2`$. The phonon transition rate due to piezoelectric coupling to the acoustic phonon is
$`\mathrm{\Gamma }_{l_1l_1^{}l_2l_2^{}}^{PZ}`$ $`=\delta _{m_1m_2,m_1^{}m_2^{}}\sqrt{{\displaystyle \frac{(n_{1+,<})!(n_{1,<})!(n_{2+,<})!(n_{2,<})!}{(n_{1+,>})!(n_{1,>})!(n_{2+,>})!(n_{2,>})!}}}(1)^{\left|n_{2+}n_{2+}^{}\right|+\left|n_2n_2^{}\right|+\frac{N}{2}}{\displaystyle \frac{(eh_{14})^2}{4\pi \mathrm{}\rho }}`$ (15)
$`{\displaystyle \underset{\nu }{}}({\displaystyle \frac{l_0q_\nu }{2}})^N{\displaystyle \frac{q_\nu }{C_\nu ^2}}{\displaystyle _0^\pi }A_\nu (\theta )\left|Z(q_\nu \mathrm{cos}\theta )\right|^2e^{\frac{1}{2}(l_0q_\nu \mathrm{sin}\theta )^2}({\displaystyle \frac{l_0^2q_\nu ^2\mathrm{sin}^2\theta }{4}})(\mathrm{sin}\theta )^{N+1}๐\theta .`$
The contribution from the deformation potential interaction is
$`\mathrm{\Gamma }_{l_1l_1^{}l_2l_2^{}}^{DP}`$ $`=\delta _{m_1m_2,m_1^{}m_2^{}}\sqrt{{\displaystyle \frac{(n_{1+,<})!(n_{1,<})!(n_{2+,<})!(n_{2,<})!}{(n_{1+,>})!(n_{1,>})!(n_{2+,>})!(n_{2,>})!}}}(1)^{\left|n_{2+}n_{2+}^{}\right|+\left|n_2n_2^{}\right|+\frac{N}{2}}{\displaystyle \frac{\mathrm{\Xi }_d^2q_l^3}{4\pi \mathrm{}\rho c_l^2}}({\displaystyle \frac{q_ll_0}{2}})^N`$ (16)
$`{\displaystyle _0^\pi }(\mathrm{sin}\theta )^{N+1}\left|Z(q_l\mathrm{cos}\theta )\right|^2\mathrm{exp}({\displaystyle \frac{q_l^2l_0^2\mathrm{sin}^2\theta }{2}})({\displaystyle \frac{q_l^2l_0^2\mathrm{sin}^2\theta }{4}})๐\theta .`$
In the above, we have used the short notation $`l`$ to denote the quantum number $`(n,m)`$ and $`n_{j,<}(n_{j,>})`$ to denote $`\mathrm{min}\{n_j,n_j^{}\}`$ ($`\mathrm{max}\{n_j,n_j^{}\}`$), e.g., $`l_1`$ stands for $`(n_1,m_1)`$ and $`n_{1+,>}`$ stands for $`\mathrm{max}\{n_{1+},n_{1+}^{}\}`$. Other quantities are $`N=\left|n_{1+}n_{1+}^{}\right|+\left|n_1n_1^{}\right|+\left|n_{2+}n_{2+}^{}\right|+\left|n_2n_2^{}\right|,`$ $`q_\nu =\left|\mathrm{\Delta }_{ll^{}}\right|/(\mathrm{}C_\nu ),`$ $`C_\nu `$ ($`\nu =l,t`$) is the longitudinal or transverse sound velocity, $`\theta `$ is the polar angle of the phonon vector $`๐ช,`$ $`A_\nu (\theta )`$ is the anisotropy function of the piezoelectric interaction, with the dependence on the the azimuth angle $`\phi `$ averaged out,
$$Z(q)=e^{iqz}\left|\chi (z)\right|^2๐z=\frac{4\pi ^2i(e^{iqz_0}1)}{qz_0\left[(qz_0)^2(2\pi )^2\right]}$$
(17)
is the form factor, $`(x)=L_{n_{1+,<}}^{\left|n_{1+}n_{1+}^{}\right|}(x)L_{n_{1,<}}^{\left|n_1n_1^{}\right|}(x)L_{n_{2+,<}}^{\left|n_{2+}n_{2+}^{}\right|}(x)L_{n_{2,<}}^{\left|n_2n_2^{}\right|}(x),`$ and $`\mathrm{\Xi }_d`$ is the deformation potential constant.
For the spin-flip transitions between the lowest Zeeman doublet $`|00+`$ and $`|00,`$ the following overlap integral and phonon transition rates are used:
$$S_{00,n_1m_1,00,n_2m_2}=\delta _{m_1,m_2}\frac{(n_1+n_2+\left|m_1\right|)!}{2^{n_1+n_2+\left|m_1\right|}\sqrt{n_1!n_2!(n_1+\left|m_1\right|)!(n_2+\left|m_2\right|)!}}\frac{\mathrm{\Omega }\xi }{2\pi z_0l_0^2},$$
(18)
$`\mathrm{\Gamma }_{00,n_1m_1,00,n_2m_2}^{PZ}`$ $`=\delta _{m_1,m_2}{\displaystyle \frac{(1)^{n_1+n_2}}{\sqrt{n_1!(n_1+\left|m_1\right|)!n_2!(n_2+\left|m_2\right|)!}}}{\displaystyle \frac{(eh_{14})^2}{4\pi \mathrm{}\rho }}{\displaystyle \underset{\nu }{}}({\displaystyle \frac{l_0q_\nu }{2}})^{2(n_1+n_2+\left|m_1\right|)}{\displaystyle \frac{q_\nu }{C_\nu ^2}}`$ (19)
$`{\displaystyle _0^\pi }A_\nu (\theta )\left|Z(q_\nu \mathrm{cos}\theta )\right|^2\mathrm{exp}({\displaystyle \frac{l_0^2q_\nu ^2\mathrm{sin}^2\theta }{2}})({\displaystyle \frac{l_0^2q_\nu ^2\mathrm{sin}^2\theta }{4}})(\mathrm{sin}\theta )^{2(n_1+n_2+\left|m_1\right|)+1}๐\theta ,`$
$`\mathrm{\Gamma }_{00,n_1m_1,00,n_2m_2}^{DP}`$ $`=\delta _{m_1,m_2}{\displaystyle \frac{(1)^{n_1+n_2}}{\sqrt{n_1!n_2!(n_1+\left|m_1\right|)!(n_2+\left|m_2\right|)!}}}{\displaystyle \frac{\mathrm{\Xi }_d^2q_l^3}{4\pi \mathrm{}\rho c_l^2}}({\displaystyle \frac{q_ll_0}{2}})^{2(n_1+n_2+\left|m_1\right|)}`$ (20)
$`{\displaystyle _0^\pi }(\mathrm{sin}\theta )^{2(n_1+n_2+\left|m_1\right|)+1}\left|Z(q_l\mathrm{cos}\theta )\right|^2\mathrm{exp}({\displaystyle \frac{q_l^2l_0^2\mathrm{sin}^2\theta }{2}})({\displaystyle \frac{q_l^2l_0^2\mathrm{sin}^2\theta }{4}})๐\theta ,`$
and $`S_{n_1m_1,n_2m_2,00,00}`$, $`\mathrm{\Gamma }_{00,n_1m_1,n_2m_2,00}^{PZ}`$, $`\mathrm{\Gamma }_{00,n_1m_1,n_2m_2,00}^{DP}`$ can be obtained from $`S_{00,n_1m_1,00,n_2m_2},`$ $`\mathrm{\Gamma }_{00,n_1m_1,00,n_2m_2}^{PZ},`$ $`\mathrm{\Gamma }_{00,n_1m_1,00,n_2m_2}^{DP}`$, respectively, by replacing $`\delta _{m_1,m_2}`$ with $`\delta _{m_1,m_2}`$.
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# On varieties of almost minimal degree in small codimension
## 1. Introduction
Let $`X_K^r`$ denote an irreducible and reduced projective variety over an algebraically closed field $`K.`$ We always assume that $`X`$ is non-degenerate, that is not contained in a hyperplane. Then, the degree and the codimension of $`X`$ satisfy the inequality $`\mathrm{deg}X\mathrm{codim}X+1`$ (cf. for instance ). Varieties for which equality holds are called called varieties of minimal degree. These varieties are completely classified (cf. for instance and ). In particular they are arithmetically Cohen-Macaulay and have a linear minimal free resolution. In particular, their Betti numbers are explicitly known.
In case $`\mathrm{deg}X=\mathrm{codim}X+2,`$ the variety $`X`$ is called a variety of almost minimal degree. Here one has a much greater variety of possible Betti numbers. The investigation of homological properties of varieties of almost minimal degree was initiated by Hoa, Stรผckrad, and Vogel (cf. ). We refer also to and for certain improvements of their results. In particular the Castelnuovo-Mumford regularity of a variety $`X`$ of almost minimal degree satisfies $`\mathrm{reg}X2`$ (cf. for the definition of the Castelnuovo-Mumford regularity).
In the framework of polarized varieties of $`\mathrm{\Delta }`$-genus 1, Fujita (cf. and ) provides a satisfactory description of varieties of almost minimal degree. The study of varieties of almost minimal degree is pursued by the authors (cf. ) from the arithmetic point of view. It turns out that $`X_K^r`$ is either an arithmetically normal Del Pezzo variety or a proper projection of a variety of minimal degree. By a proper projection of a variety $`Z_K^{r+1}`$ we always mean a projection from a point $`P_K^{r+1}Z.`$ See also 3.1 for the precise statement.
The aim of the present paper is to investigate varieties of almost minimal degree and of low codimension, in particular their Betti diagrams. More precisely, we describe the structure of the minimal free resolution of a variety $`X`$ of almost minimal degree of $`\mathrm{codim}X4`$ by listing all possible Betti diagrams. Let us recall that the structure of arithmetically Cohen-Macaulay resp. Gorenstein varieties in codimension 2 resp. 3 is known by the Theorems of Hilbert-Burch resp. Buchsbaum-Eisenbud (cf. ). So, we need not to discuss these cases in detail any more. As the Betti diagram, the degree and the codimension are not affected if $`X`$ is replaced by a cone over $`X`$, we shall assume that $`X`$ is not a cone. The most surprising fact is, that the dimension of $`X`$ is always $`6`$ (cf. Section 2 for the precise statements). Our main technical tool is a result shown by the authors in , which says that apart from an exceptional case, (that is the generic projection of the Veronese surface in $`_K^5`$) any non-arithmetically normal (and in particular non-arithmetically Cohen-Macaulay) variety of almost minimal degree $`X_K^r`$ (which is not a cone) is contained in a variety of minimal degree $`Y_K^r`$ such that $`\mathrm{codim}(X,Y)=1.`$
###### Acknowledgement .
The authors thank Uwe Nagel for making his preprint available. They also thank for the refereeโs very valuable suggestions.
## 2. Main Results
Let $`X_K^r`$ denote a non-degenerate variety of almost minimal degree, hence an integral closed subscheme with $`\mathrm{deg}X=\mathrm{codim}X+2`$ not contained in a hyperplane $`_K^{r1}_K^r.`$ We use the abbreviations $`dimX=d`$ and $`\mathrm{codim}X=c.`$ Let $`S=K[x_0,\mathrm{},x_r]`$ denote the polynomial ring in $`r+1`$ variables, so that $`_K^r=\mathrm{Proj}(S).`$ Let $`A_X=S/I_X`$ denote the homogeneous coordinate ring of $`X,`$ where $`I_XS`$ is the defining ideal of $`X.`$ The $`\mathrm{codepth}`$ of $`A_X`$ is defined as the difference $`\mathrm{codepth}A_X:=dimA_X0ptA_X,`$ where $`0ptA_X`$ denotes the depth of $`A_X.`$
For the notion of Betti diagrams we follow the suggestion of Eisenbud (cf. ). That is, in a diagram the number in the $`i`$-th column and the $`j`$-th row is
$$dim_K\mathrm{Tor}_i^S(K,S/I_X)_{i+j},i1.$$
Outside the range of the diagram all the entries are understood to be zero. Our main results are.
###### Theorem 2.1.
$`(\mathrm{codim}X=2)`$ Let $`X_K^r`$ be a non-degenerate variety of degree 4 and codimension 2 which is not a cone. Then $`X`$ is of one of the following types:
* $`X`$ is a complete intersection cut out by two quadrics.
* $`dimX4`$ and the Betti diagram of $`I_X`$ has the form
| | 1 | 2 | 3 |
| --- | --- | --- | --- |
| 1 | 1 | 0 | 0 |
| 2 | 3 | 4 | 1 |
.
* (The exceptional case) $`X`$ is a generic projection of the Veronese surface $`F_K^5`$ and the Betti diagram of $`I_X`$ has the shape
| | 1 | 2 | 3 | 4 |
| --- | --- | --- | --- | --- |
| 1 | 0 | 0 | 0 | 0 |
| 2 | 7 | 10 | 5 | 1 |
.
Moreover for any $`1d4`$ there are examples as mentioned in (b) such that $`dimX=d.`$
In the case where $`X`$ is a Cohen-Macaulay variety, Theorem 2.1 (b) has been shown by Nagel (cf. ). Under this additional assumption one has $`dimX2.`$ Theorem 2.1 grew out of our aim to understand Nagelโs arguments. Our approach enables us to investigate the cases of codimension three and four as well.
###### Theorem 2.2.
$`(\mathrm{codim}X=3)`$ Let $`X_K^r`$ denote a non-degenerate variety of degree 5 and codimension 3 which is not a cone. Then the following cases may occur:
* $`dimX6`$ and $`X`$ is the Pfaffian variety defined by the five Pfaffians of a skew symmetric $`5\times 5`$ matrix of linear forms. Its minimal free resolution is given by the Buchsbaum-Eisenbud complex.
* $`dimX4`$ and $`\mathrm{codepth}A_X=1.`$ The Betti diagram of the defining ideal $`I_X`$ has the following form
| | 1 | 2 | 3 | 4 |
| --- | --- | --- | --- | --- |
| 1 | 4 | 2 | 0 | 0 |
| 2 | 1 | 6 | 5 | 1 |
.
* $`dimX5`$ and $`\mathrm{codepth}A_X=2.`$ The Betti diagram of $`I_X`$ has the following shape
| | 1 | 2 | 3 | 4 | 5 |
| --- | --- | --- | --- | --- | --- |
| 1 | 3 | 2 | 0 | 0 | 0 |
| 2 | 6 | 16 | 15 | 6 | 1 |
.
For any of the dimensions and codepths admitted in (a), (b) and (c) resp. there are examples of varieties of almost minimal degree.
The next result concerns the case where $`X`$ is of codimension 4. A new phenomenon occurs in this situation. Namely, for the same codepth two different Betti diagrams may occur.
###### Theorem 2.3.
$`(\mathrm{codim}X=4)`$ Let $`X_K^r`$ denote a non-degenerate variety of degree 6 and codimension 4 which is not a cone. Then the following four cases may occur:
* $`dimX4`$ and $`X`$ is arithmetically Gorenstein. Its minimal free resolution has the following form
$$0S(6)S^9(4)S^{16}(3)S^9(2)I_X0.$$
* $`dimX4`$ and $`\mathrm{codepth}A_X=1.`$ The Betti diagram of the defining ideal $`I_X`$ has one of the following two forms:
| | 1 | 2 | 3 | 4 | 5 |
| --- | --- | --- | --- | --- | --- |
| 1 | 8 | 12 | 3 | 0 | 0 |
| 2 | 1 | 4 | 10 | 6 | 1 |
resp. 1 2 3 4 5 1 8 11 3 0 0 2 0 4 10 6 1 .
* $`dimX5`$ and $`\mathrm{codepth}A_X=2.`$ The Betti diagram of $`I_X`$ is of the form:
| | 1 | 2 | 3 | 4 | 5 | 6 |
| --- | --- | --- | --- | --- | --- | --- |
| 1 | 7 | 8 | 3 | 0 | 0 | 0 |
| 2 | 3 | 19 | 30 | 21 | 7 | 1 |
.
* $`dimX6`$ and $`\mathrm{codepth}A_X=3.`$ The Betti diagram of $`I_X`$ has the following shape:
| | 1 | 2 | 3 | 4 | 5 | 6 | 7 |
| --- | --- | --- | --- | --- | --- | --- | --- |
| 1 | 6 | 8 | 3 | 0 | 0 | 0 | 0 |
| 2 | 10 | 40 | 65 | 56 | 28 | 8 | 1 |
.
For any of the dimensions, codepths and Betti diagrams admitted in (a), (b), (c) and (d) resp. there are examples of varieties of almost minimal degree.
The final result of this note concerns varieties of almost minimal degree and of codimension less than half the embedding dimension. First of all note (cf. Example 4.8) that in codimension 6 there are varieties of almost minimal degree having the same $`\mathrm{codepth}`$ but with rather different Betti diagrams. That is, the corresponding statements to Theorem 2.1, 2.2 and 2.3 are not true in hihgher codimension.
###### Corollary 2.4.
Let $`X_K^r`$ be a non-degenerate variety of almost minimal degree which is not a cone. Suppose that $`dimX>\mathrm{codim}X+2`$ and $`\mathrm{codim}X3.`$ Then $`\mathrm{codim}X=3`$ and $`X`$ is arithmetically Gorenstein. Therefore $`dimX6`$ and $`X`$ is defined by the five Pfaffians of a skew symmetric $`5\times 5`$ matrix of linear forms.
Moreover, in the case where $`X`$ is a variety of almost minimal degree and $`dimX\mathrm{codim}X+2,`$ that is when $`X`$ is not necessary a Del Pezzo variety there are estimates for the Betti numbers and their vanishing (cf. 3.3 for the details).
## 3. Outline of the Proofs
Let $`X_K^r`$ denote a non-degenerate reduced irreducible variety of almost minimal degree. By the work of the authors (cf. ) it follows that $`X`$ is either a normal Del Pezzo variety โ in this case $`X`$ is arithmetically Gorenstein โ or a projection of a variety of minimal degree.
First, we consider the case of non-arithmetically normal varieties of almost minimal degree. In this situation we have the following characterization in which $`\mathrm{Sec}_P(Z)`$ denotes the secant cone of a projective variety $`Z`$ with respect to a point $`P`$ in the ambient space.
###### Lemma 3.1.
Let $`X_K^r`$ denote a non-degenerate reduced irreducible variety which is not a cone. Let $`1tdimX+1=:d+1.`$ Then the following conditions are equivalent:
* $`X`$ is not arithmetically normal, $`\mathrm{deg}X=\mathrm{codim}X+2`$ and $`0ptA_X=t.`$
* $`X`$ is the projection of a variety $`Z_K^{r+1}`$ of minimal degree from a point $`P_K^{r+1}Z`$ such that $`dim\mathrm{Sec}_P(Z)=t1.`$
Moreover, $`10ptA_X4.`$
###### Proof.
Cf. \[2, Theorem 1.1 and Corollary 7.6\]. โ
In view of Lemma 3.1 there is some need for information about varieties of minimal degree in order to understand varieties of almost minimal degree. A variety of minimal degree $`Z_K^s`$ is either
* a quadric hypersurface,
* a (cone over a) Veronese surface in $`_K^5,`$ or
* a (cone over a) smooth rational normal scroll
(cf. and for the details and the history of this classification).
Next we recall a few basic facts about rational normal scrolls (cf. also ). Let
$$T=K[x_{10},\mathrm{},x_{1a_1},x_{20},\mathrm{},x_{2a_2},\mathrm{},x_{k0},\mathrm{},x_{ka_k}]$$
be the polynomial ring. A (cone over a) rational normal scroll $`S(a_1,\mathrm{},a_k)`$ is defined as the โrank two subscheme in $`_K^s=\mathrm{Proj}(T)`$โ defined by the matrix
$$M=\left(\begin{array}{ccccccccccccc}x_{10}& \mathrm{}& x_{1a_11}& \mathrm{}& x_{20}& \mathrm{}& x_{2a_21}& \mathrm{}& \mathrm{}& \mathrm{}& x_{k0}& \mathrm{}& x_{ka_k1}\\ x_{11}& \mathrm{}& x_{1a_1}& \mathrm{}& x_{21}& \mathrm{}& x_{2a_2}& \mathrm{}& \mathrm{}& \mathrm{}& x_{k1}& \mathrm{}& x_{ka_k}\end{array}\right)$$
with $`s=k1+_{i=1}^ka_i.`$ It is well known that $`dimS(a_1,\mathrm{},a_k)=k`$ and therefore
$$\mathrm{deg}S(a_1,\mathrm{},a_k)=\mathrm{codim}S(a_1,\mathrm{},a_k)+1=\underset{i=1}{\overset{k}{}}a_i,$$
so that $`S(a_1,\mathrm{},a_k)`$ is a variety of minimal degree. Keep in mind that $`S(a_1,\mathrm{},a_k)`$ is a proper cone if and only if $`a_i=0`$ for some $`i\{1,\mathrm{},k\},`$ that is, if and only if there are indeterminates which do not occur in $`M.`$
Moreover we need some information about the Hilbert series $`F(\lambda ,A_X)`$ of the graded $`K`$-algebra $`A_X.`$ Let us recall that the Hilbert series of a graded $`K`$-algebra $`A`$ is the formal power series defined by
$$F(\lambda ,A)=\underset{i0}{}(dim_KA_i)\lambda ^i.$$
The Hilbert series of a variety of almost minimal degree may be described as follows.
###### Lemma 3.2.
Let $`X_K^r`$ denote a variety of almost minimal degree. Put $`q=\mathrm{codepth}A_X.`$ Then
$$F(\lambda ,A_X)=\frac{1}{(1\lambda )^{d+1}}(1+(c+1)\lambda \lambda (1\lambda )^{q+1}),$$
where $`c=\mathrm{codim}X`$ and $`d=dimX.`$ Furthermore $`dim_K(I_X)_2=\left(\genfrac{}{}{0pt}{}{c+1}{2}\right)q1.`$
###### Proof.
Cf. \[2, Corollary 4.4\]. โ
As a consequence of Lemma 3.1 the authors (cf. ) derived some information about the Betti numbers of $`I_X`$ for certain varieties $`X_K^r`$ of almost minimal degree.
###### Lemma 3.3.
Let $`X_K^r`$ be a variety of almost minimal degree which is not arithmetically Cohen-Macaulay. Suppose that $`X`$ is not a generic projection of (a cone over) the Veronese surface $`F_K^5.`$ Then there exists a variety of minimal degree $`Y_K^r`$ such that $`XY`$ and $`\mathrm{codim}(X,Y)=1.`$ Moreover
$$\mathrm{Tor}_i^S(k,A_X)k^{u_i}(i1)k^{v_i}(i2)\text{ for }1ic+q,$$
where $`c=\mathrm{codim}X,q=\mathrm{codepth}A_X`$ and
* + $`u_1=\left(\genfrac{}{}{0pt}{}{c+1}{2}\right)q1,`$
+ $`i\left(\genfrac{}{}{0pt}{}{c}{i+1}\right)u_i(c+1)\left(\genfrac{}{}{0pt}{}{c}{i}\right)\left(\genfrac{}{}{0pt}{}{c}{i+1}\right),\text{if}1<i<cq,`$
+ $`u_i=i\left(\genfrac{}{}{0pt}{}{c}{i+1}\right),\text{if}cqi<c,`$
+ $`u_i=0,\text{if}cic+q.`$
* + $`\mathrm{max}\{0,\left(\genfrac{}{}{0pt}{}{c+q1}{i+1}\right)(i+2)\left(\genfrac{}{}{0pt}{}{c}{i+1}\right)\}v_i\left(\genfrac{}{}{0pt}{}{c+q+1}{i+1}\right),\text{if}1i<cq1,`$
+ $`v_i=\left(\genfrac{}{}{0pt}{}{c+q+1}{i+1}\right)(i+2)\left(\genfrac{}{}{0pt}{}{c}{i+1}\right),\text{if}\mathrm{max}\{1,cq1\}i<c,`$
+ $`v_i=\left(\genfrac{}{}{0pt}{}{c+q+1}{i+1}\right),\text{if}cic+q.`$
In addition $`v_iu_{i+1}=\left(\genfrac{}{}{0pt}{}{c+q+1}{i+1}\right)(c+1)\left(\genfrac{}{}{0pt}{}{c}{i+1}\right)+\left(\genfrac{}{}{0pt}{}{c}{i+2}\right)`$ for all $`1i<c.`$
###### Proof.
Cf. \[2, Theorem 1.1 and Theorem 8.3\]. โ
In the particular case where $`dimX=1`$ the statement about the Betti numbers has been shown independently by Nagel (cf. ).
In the following remark, we add a comment concerning the โexceptional casesโ of a generic projection of (a cone over) the Veronese surface $`F_K^5`$ and of an arithmetically Cohen-Macaulay variety.
###### Remark 3.4.
A) (The exceptional case) Let $`F_K^5`$ be the Veronese surface defined by the $`2\times 2`$-minors of the symmetric matrix
$$M=\left(\begin{array}{ccc}x_0& x_1& x_2\\ x_1& x_3& x_4\\ x_2& x_4& x_5\end{array}\right).$$
Let $`P_K^5F`$ denote a point. Suppose that $`\mathrm{rank}M_P=3,`$ that is the case of a generic point. Remember that $`detM=0`$ defines the secant variety of $`F.`$ Then the projection of $`F`$ from $`P`$ defines a surface $`X_K^4`$ of almost minimal degree and $`0ptA_X=1.`$ The surface $`X`$ is cut out by seven cubics (cf. 3.2), so that it is not contained in a variety of minimal degree.
B) (The arithmetically normal case) Let $`X_K^r`$ denote an arithmetically normal variety of almost minimal degree. Then $`X`$ is not a birational projection of a scroll and hence a maximal Del Pezzo variety (cf. \[2, Theorem 1.2\]). In particular $`X`$ is arithmetically Gorenstein. If in addition $`\mathrm{codim}X3,`$ Fujitaโs classification of normal maximal Del Pezzo varieties yields $`dimX4`$ (cf. \[5, (8.11), (9.17)\]). If $`\mathrm{codim}X4,`$ the same classification shows that $`dimX4.`$
C) If $`X`$ is an arithmetical normal Del Pezzo variety it is in general not a one codimensional subvariety of a variety of minimal degree. Namely, let $`X_K^9`$ be the smooth codimension three variety cut out by the $`4\times 4`$-Pfaffians of a generic skew-symmetric $`5\times 5`$-matrix of linear forms. Then $`X`$ is not contained in a variety $`Y`$ of minimal degree such that $`\mathrm{codim}(X,Y)=1`$ (cf. for the details).
Proofs. After these preparations, we now come to the proofs of our Theorems. The proof of statement 2.1 (a) is easy. Now let us consider the statements 2.2 (a) and 2.3 (a). In both cases $`X`$ is arithmetically Gorenstein. If $`X`$ is a birational projection of a scroll, we have $`dimA_X=0ptA_X4`$ (cf. Lemma 3.1). Otherwise $`X`$ is arithmetically normal (cf. \[2, Theorem 1.2\]). So, by 3.4 B) we have $`dimX6`$ if $`\mathrm{codim}X3`$ and $`dimX4`$ if $`\mathrm{codim}X4.`$ This gives us the dimension estimates. The statements on the minimal free resolutions now follow from well known results. In fact the structure of the minimal free resolution of $`I_X`$ given in Theorem 2.3 (a) is a consequence of \[13, Theorem B\].
Next we prove statement (c) of Theorem 2.1. To this end let $`X_K^4`$ be a generic projection of the Veronese surface $`F_K^5.`$ Then $`dimX=2,\mathrm{codim}X=2`$ and $`0ptA_X=1.`$ As seen above, $`I_X`$ does not contain any quadric. Therefore $`I_X`$ has a linear resolution. Remember that $`\mathrm{reg}I_X=3`$ (cf. ). A computation with the aid of the Hilbert series (cf. 3.2) gives the structure of the Betti diagram of 2.1 (c).
For all other statements of Theorems 2.1, 2.2 and 2.3 we may assume that the variety of almost minimal degree $`X_K^r`$ is not arithmetically Cohen-Macaulay and not a projection of the Veronese surface $`F_K^5.`$ So, $`X`$ is contained in a variety of minimal degree $`Y`$ of codimension $`c1`$ (cf. 3.3). That is, $`Y`$ is defined as the zero locus of $`\left(\genfrac{}{}{0pt}{}{c}{2}\right)`$ quadrics. On the other hand the defining ideal $`I_X`$ is generated by $`\left(\genfrac{}{}{0pt}{}{c+1}{2}\right)q1`$ quadrics (cf. 3.2). This implies that
$$dim_K(I_X/I_Y)_2=cq10,$$
where $`c=\mathrm{codim}X,q=\mathrm{codepth}A_X1.`$ Considering all possibilities that arise for $`c=2,3,4`$ it follows that, with the exception of the case in which $`c=4,q=1,`$ Lemma 3.3 furnishes the corresponding Betti diagrams.
The particular case where $`c=4,q=1,`$ yields the following shape of the Betti diagram
| | 1 | 2 | 3 | 4 | 5 |
| --- | --- | --- | --- | --- | --- |
| 1 | 8 | $`u_2`$ | 3 | 0 | 0 |
| 2 | $`v_1`$ | 4 | 10 | 6 | 1 |
with $`v_1u_2=11.`$ In order to finish the proof we observe that $`v_11`$ (cf. Lemma 3.5).
To complete the proof of Theorems 2.1, 2.2 and 2.3 we have to prove the stated constraints on the occuring dimensions and codepths. To do so, we may assume that $`X`$ is not arithmetically Cohen-Macaulay. Therefore by Lemma 3.1 $`X`$ is the projection of a variety of minimal degree $`Z_K^{r+1}`$ from a point $`P_K^{r+1}Z`$ such that $`dim\mathrm{Sec}_P(Z)=t1,`$ where $`t=0ptA_X.`$
Next let us analyze this situation in more detail. To this end let $`Z=S(a_1,\mathrm{},a_k)`$ for certain integers $`a_i,i=1,\mathrm{},k.`$ Then it follows that
$$r=k2+\underset{i=1}{\overset{k}{}}a_i\text{ and }c+2=\underset{i=1}{\overset{k}{}}a_i,$$
where $`c=\mathrm{codim}X.`$
As $`X`$ is not a cone, $`Z`$ cannot be a cone over a rational normal scroll. Therefore $`\mathrm{min}\{a_i:i=1,\mathrm{},k\}1.`$ So, for a given codimension $`c`$ we have to investigate all the possible partitions
$$c+2=\underset{i=1}{\overset{k}{}}a_i,\text{ with }k1\text{ and }a_1a_2\mathrm{}a_k1.$$
For $`c=2`$ we thus get the following possible types for the rational normal scroll $`Z:`$
$$\begin{array}{ccccccc}k& a_1& a_2& a_3& a_4& r+1& dimX\\ & & & & & & \\ 1& 4& & & & 4& 1\\ 2& 3& 1& & & 5& 2\\ 2& 2& 2& & & 5& 2\\ 3& 2& 1& 1& & 6& 3\\ 4& 1& 1& 1& 1& 7& 4\end{array}$$
This proves already that $`dimX4.`$
Next, we discuss the case in which the codimension equals $`3.`$ Here, there are the following possibilities for the type of the scroll $`Z`$:
$$\begin{array}{cccccccc}k& a_1& a_2& a_3& a_4& a_5& r+1& dimX\\ & & & & & & & \\ 1& 5& & & & & 5& 1\\ 2& 4& 1& & & & 6& 2\\ 2& 3& 2& & & & 6& 2\\ 3& 3& 1& 1& & & 7& 3\\ 3& 2& 2& 1& & & 7& 3\\ 4& 2& 1& 1& 1& & 8& 4\\ 5& 1& 1& 1& 1& 1& 9& 5\end{array}$$
Therefore $`dimX5.`$ We already know that $`u_13.`$ Now, on use of Lemma 3.3 we easily get the requested constraints in Theorem 2.2.
Finally, let $`\mathrm{codim}X=4.`$ Then, as above there is the following list of possible types for the scroll $`Z:`$
$$\begin{array}{ccccccccc}k& a_1& a_2& a_3& a_4& a_5& a_6& r+1& dimX\\ & & & & & & & & \\ 1& 6& & & & & & 6& 1\\ 2& 5& 1& & & & & 7& 2\\ 2& 4& 2& & & & & 7& 2\\ 2& 3& 3& & & & & 7& 2\\ 3& 4& 1& 1& & & & 8& 3\\ 3& 3& 2& 1& & & & 8& 3\\ 3& 2& 2& 2& & & & 8& 3\\ 4& 3& 1& 1& 1& & & 9& 4\\ 4& 2& 2& 1& 1& & & 9& 4\\ 5& 2& 1& 1& 1& 1& & 10& 5\\ 6& 1& 1& 1& 1& 1& 1& 11& 6\end{array}$$
As above it follows that $`dimX6.`$ We know that $`u_16.`$ So, by Lemma 3.3 we get the requested constraints in Theorem 2.3.
Finally we prove Corollary 2.4. Assume that $`X`$ is not arithmetically normal. Then, by Theorem 1.2 of we know that $`X`$ is a birational projection of a rational normal scroll $`Z_K^{r+1}`$ from a point $`p_K^{r+1}Z.`$ As $`X`$ is not a cone, $`Z`$ is not a cone either and therefore $`\mathrm{codim}X+2=_{i=1}^ka_ik=dimX.`$ This contradicts the assumption of Corollary 2.4. Therefore $`X`$ is arithmetically normal and so $`dimX6`$ by Remark 3.4. That is, $`\mathrm{codim}X=3.`$ So our claim follows by Theorem 2.2.
For the existence of the samples described in Theorems 2.1 and 2.2 we refer to the next section.
We close this section with a result on the number of cubics in a minimal generating set of the defining ideal of a certain varieties of almost minimal degree.
###### Lemma 3.5.
Let $`X_K^r`$ be a variety of almost minimal degree with $`\mathrm{codepth}A_X=1`$ and $`c:=\mathrm{codim}X4.`$ Then the defining ideal $`I_X`$ of $`X`$ is generated by $`\left(\genfrac{}{}{0pt}{}{c+1}{2}\right)2`$ quadrics and at most one cubic.
###### Proof.
First we reduce the problem to the case in which $`dimX=1.`$ Let $`d=dimX>1.`$ By an argument of Bertini type (cf. ) we may find generic linear forms $`l_1,\mathrm{},l_{d1}S_1`$ such that $`W=X_K^{c+1}_K^{c+1}:=\mathrm{Proj}(S/(l_0,\mathrm{},l_{d1})S)`$ is a non-degenerate integral variety of almost minimal degree. As $`l_0,\mathrm{},l_{d1}`$ are chosen generically and $`0ptA_X=d`$ they form an $`A_X`$-regular sequence. Therefore the Betti diagrams of $`I_X`$ and $`I_W`$ are the same. In particular $`\mathrm{codepth}A_W=1.`$
So we assume $`X_K^s`$ with $`dimX=1`$ and $`s=c+1.`$ The statement about the number of quadrics is a consequence of Lemma 3.2. Since $`X`$ is of almost minimal degree we know that $`I=I_X`$ is 3-regular (cf. 3.3). Write $`I=(J,LS)`$ with $`J=I_2S`$ and with a $`K`$-vector space $`LS_3`$ such that $`I_3=J_3L.`$ Our aim is to show that $`dim_KL1.`$
After an appropriate linear coordinate change we may assume that $`x_sS_1`$ is generic. Let $`T:=S/x_sS=K[x_0,\mathrm{},x_{s1}].`$ Then $`R:=T/(J,L)TS/(I,x_sS)`$ defines a scheme $`Z`$ of $`s+1`$ points in semi uniform position in $`_K^{s1}.`$ The short exact sequence $`0A_X(1)\stackrel{x_s}{}A_XR0`$ induces an isomorphism $`H_{R_+}^0(R)K(2).`$ But this means that the vanishing ideal of $`Z`$ in $`T`$ has the form $`(J,L,q)T`$ with an appropriate quadric $`qS_2.`$ Since $`s3`$ the minimal free resolution of this ideal has the form
$$T^{a_2}(3)\stackrel{\varphi }{}T^{a_1}(2)\stackrel{\pi }{}(J,L,q)T0.$$
This allows us to write $`(J,L,q)T=(J,q)T`$ and to assume that the first $`a_11`$ generators of $`T^{a_1}(2)`$ are mapped by $`\pi `$ onto a $`K`$-basis of $`(JT)_2`$ and the last generator is mapped by $`\pi `$ to $`q1_T.`$ Clearly, $`\varphi `$ is given by a matrix with linear entries. This shows that $`M:=JT:_TqT`$ is a proper ideal generated by linear forms.
As $`JTM`$ and as $`(J,qT)=IT`$ is of height $`s1`$ we must have $`s20ptMs.`$ As $`M`$ is generated by linear forms, $`(T/M)_1`$ is a $`K`$-vector space of dimension $`t\{0,1,2\}.`$ So the graded short exact sequence
$$0T/M(2)T/JTT/(J,q)T0$$
shows that
$$dim_K(IT)_3=dim_K((J,q)T)_3=dim_K(JT)_3+t.$$
Therefore, we may write $`(I,x_s)=(J,L^{},x_s)`$ where $`L^{}LS_3`$ is a $`K`$-vector space of dimension $`t.`$ As $`I`$ is a prime ideal and as $`x_sS_1I`$ it follows $`I=(J,L^{}),`$ hence $`L^{}=L.`$ So, if $`dim_KL^{}1,`$ we are done.
Otherwise, $`dim_KL^{}=dim_KL=2=t`$ and we may write $`I=(J,k_1,k_2)`$ with $`k_1,k_2S_3.`$ As $`0ptI=s1`$ it follows $`0ptJs3.`$ As $`JTM`$ and as $`0ptM=s2`$ we have $`0ptJTs2.`$ As $`x_s`$ is a generic linear form, this means that $`0ptJs3`$ and hence $`0ptJ=s3.`$ As $`I=(J,k_1,k_2)`$ is a prime ideal of height $`s1=0ptJ+2,`$ the ideal $`J`$ must be prime too.
As $`x_s`$ is generic and $`0ptJs3,`$ we may conclude by Bertiniโs theorem that $`JTT`$ defines an integral subscheme of $`_K^{s1}=\mathrm{Proj}(T).`$ So, the saturation
$$JT:_TT_+T\text{ of }JT\text{ in }T$$
is a prime ideal of height $`s2.`$ As $`JTMT_+`$ and as $`M`$ is a prime ideal we get $`JT:_TT_+=M.`$ Therefore
$$\mathrm{Proj}(T/IT)=\mathrm{Proj}(T/(J,q)T)=\mathrm{Proj}(T/(M,qT))$$
consists of two points, so that $`s+1=2,`$ a contradiction. So, the case $`dim_KL^{}2`$ does not occur at all. โ
The previous Lemma 3.5 is inspired by a corresponding statement for curves of degree $`r+2`$ in $`_K^r`$ shown by the authors (cf. \[1, Lemma (6.4)\]).
Moreover, if $`\mathrm{codepth}A_X2`$ the number of cubics needed to define $`X`$ is not bounded by 1 (cf. the examples in \[2, Section 9\]).
## 4. Examples
In this section we want to confirm the existence of all types of varieties of almost minimal degree $`XP_K^r`$ which are described in the Theorems 2.1, 2.2 and 2.3.
First of all we want to show the existence of a Del Pezzo variety as required by Theorem 2.3 (a).
###### Example 4.1.
Let $`X=_K^2\times _K^2_K^8`$ be the Segre product of two projective planes. Its defining ideal $`I_X`$ is generated by the $`2\times 2`$-minors of the the following generic $`3\times 3`$-matrix
$$\left(\begin{array}{ccc}x_0& x_1& x_2\\ x_3& x_4& x_5\\ x_6& x_7& x_8\end{array}\right).$$
It is easy to see that $`dimX=4,\mathrm{codim}X=4`$ and $`\mathrm{deg}X=6.`$ Therefore, $`X`$ is a variety of almost minimal degree. Moreover, $`A_X`$ is a Cohen-Macaulay and therefore a Gorenstein ring. An example of dimension 3 is the Segre product $`_K^1\times _K^1\times _K^1_K^7`$ (cf. \[5, (8.11), 6)\]). Examples of smaller dimensions are obtained by taking generic linear sections.
In the next examples let us show that the two different Betti diagrams of statement (b) in 2.3 indeed occur. Note that they require $`\mathrm{codepth}A_X=1`$ and $`\mathrm{codim}X=4.`$
###### Example 4.2.
Consider the rational normal surface scroll $`Z=S(3,3)_K^7.`$ Let $`P_1=(0:0:0:0:1:0:0:1)`$ and $`P_2=(1:0:0:0:0:0:1)`$ in $`_K^7.`$ Then $`P_i_K^7Z,i=1,2,`$ as it is easily seen. Define $`X_i`$ to be the projection of $`Z`$ from $`P_i,i=1,2.`$ Then $`dimX_i=2`$ and $`\mathrm{codepth}A_{X_i}=1`$ for $`i=1,2.`$ The Betti diagrams of $`I_{X_i},i=1,2,`$ are those of Theorem 2.3 (b).
Now we construct non-arithmetically Cohen-Macaulay varieties of almost minimal degree of the type mentioned in Theorem 2.2. To this end we use the possible rational normal scrolls $`Z=S(a_1,\mathrm{},a_k)`$ of the proof of Theorem 2.2 which after appropriate projection furnish the varieties we are looking for. The construction of the examples corresponding to Theorems 2.1 and 2.3 follows similarly, and so we skip the details in these cases.
###### Example 4.3.
Let $`Z=S(5)_K^5`$ denote the rational normal curve of degree 5. Choose $`P_K^5Z`$ a generic point. Then, the projection $`X_K^4`$ of $`Z`$ from $`P`$ is an example of a variety of almost minimal degree with $`dimX=1`$ and $`\mathrm{codepth}A_X=1.`$
Next we want to investigate the case of surfaces.
###### Example 4.4.
Let $`Z=S(4,1)_K^6.`$ Consider the two points $`P_1=(0:1:0:0:0:0:0)`$ and $`P_2=(0:0:1:0:0:0:0).`$ Then $`P_i_K^6Z,`$ for $`i=1,2.`$ Let $`X_i,i=1,2,`$ denote the projection of $`Z`$ from $`P_i.`$ Then $`\mathrm{codepth}A_{X_1}=1`$ and $`\mathrm{codepth}A_{X_2}=2.`$ The same type of examples may be produced by projections of the scroll $`S(3,2).`$
Our next examples are of dimension 3.
###### Example 4.5.
Let $`Z=S(3,1,1)_K^7.`$ Consider the points $`P_1=(0:1:0:0:0:0:0:0)`$ and $`P_2=(0:0:0:1:1:0:0:0).`$ Then it is easy to see that $`P_i_K^7Z,i=1,2.`$ Let $`X_i`$ denote the projection of $`Z`$ from $`P_i.`$ Then $`\mathrm{codepth}A_{X_1}=2`$ and $`\mathrm{codepth}A_{X_1}=1.`$ The same type of examples may be produced by projections from the scroll $`S(2,2,1).`$
Now, let us consider the situation of fourfolds.
###### Example 4.6.
Consider the scroll $`Z=S(2,1,1,1)_K^8.`$ Let $`P_1=(0:1:0:0:0:0:0:0:0)`$ and $`P_2=(0:0:0:0:0:0:1:1:0).`$ Then $`P_i_K^8Z,i=1,2.`$ Let $`X_i^7`$ denote the projection of $`Z`$ from $`P_i,i=1,2.`$ Then $`\mathrm{codepth}A_{X_1}=2,`$ while $`\mathrm{codepth}A_{X_2}=1.`$
Finally let us consider the case where $`dimX=5.`$
###### Example 4.7.
Let $`Z=S(1,1,1,1,1)_K^9`$ be the Segre variety. Then $`P_K^9Z`$ for the point $`P=(0:1:1:0:0:0:0:0:0:0).`$ Let $`X_K^8`$ denote the projection of $`Z`$ from $`P.`$ Then $`0ptA_X=4,`$ and therefore $`\mathrm{codepth}A_X=2.`$ Finally observe that $`\mathrm{codepth}A_X=1`$ is impossible if $`dimX=5,`$ as $`0ptA_X4`$ (cf. Lemma 3.1).
The Examples 4.34.7 provide the existence of the samples claimed by Theorem 2.2. Similar constructions provide varieties as mentioned in Theorem 2.1 and 2.3.
In the final examples, we will show that in higher codimension, the shape of the Betti diagram of $`I_X`$ for a variety $`X`$ of minimal degree may vary in a much stronger way: In fact the โbeginning of the Betti diagramsโ may be rather different from each other.
###### Example 4.8.
Let $`Z=S(8)_K^8`$ denote the rational normal curve of degree 8. Let $`P_1=(0:0:0:0:0:0:1:0:0),P_2=(0:0:0:0:0:1:0:0:0)`$ and $`P_3=(0:0:0:0:1:0:0:0:0).`$ Then $`P_i_K^8Z`$ for $`i=1,2,3.`$ Let $`X_i_K^7,i=1,2,3,`$ denote the projection of $`Z`$ from $`P_i.`$ Then the Betti diagrams of $`I_{X_i},i=1,2,3,`$ resp. have the form:
| $`i`$ | | 1 | 2 | 3 | 4 | 5 | 6 | 7 |
| --- | --- | --- | --- | --- | --- | --- | --- | --- |
| 1 | 1 | 19 | 58 | 75 | 44 | 5 | 0 | 0 |
| | 2 | 1 | 6 | 15 | 20 | 21 | 8 | 1 |
| 2 | 1 | 19 | 57 | 70 | 34 | 5 | 0 | 0 |
| | 2 | 0 | 1 | 5 | 20 | 21 | 8 | 1 |
| 3 | 1 | 19 | 57 | 69 | 34 | 5 | 0 | 0 |
| | 2 | 0 | 0 | 5 | 20 | 21 | 8 | 1 |
In all three cases $`\mathrm{codim}X_i=4`$ and $`\mathrm{codepth}A_{X_i}=1.`$ Remember that the number of cubics in the defining ideals is bounded by 1 (cf. Lemma 3.5). It follows that $`X_1`$ is contained in the scroll $`S(5,1),`$ while $`X_2`$ is contained in the scroll $`S(4,2)`$ and $`X_3`$ is contained in the scroll $`S(3,3).`$
In view of the Example 4.8 and corresponding examples in higher dimensions one might expect that the type of the rational normal scroll $`Y,`$ that contains the variety $`X`$ of almost minimal degree as a one codimensional subvariety, determines the Betti diagram โnear the beginning of the resolutionโ. In small codimensions, the different types of these scrolls $`Y`$ are much more limited than imposed by Theorems 2.1, 2.2 and 2.3. It seems rather challenging to understand the rรดle of the scrolls $`Y`$ for the beginning of the minimal free resolution of $`I_X.`$
Moreover the examples in 4.8 show that the estimates for the Betti numbers given in Lemma 3.3 near the beginning of the Betti diagram are fairly weak.
###### Remark 4.9.
To compute the Betti diagrams and hence the arithmetic depths of the above examples, we have made use of the computer algebra system Singular (cf. ). Moreover, there is in preparation a conceptual approach for the computation of the $`0ptA_X`$ in terms of the center of the projection and the secant variety of $`S(a_1,\mathrm{},a_k)_K^{r+1}.`$
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# Isospin Transport at Fermi Energies
## I INTRODUCTION
In the last few years the increased accuracy of the experimental techniques has renewed interest in nuclear reactions at Fermi energies. Exclusive measurements, event-by-event analysis, and a $`4\pi `$ coverage allow a deeper investigation of the evolution of the reaction mechanisms with beam energy and centrality. New insights into the understanding of the nuclear matter equation-of-state ($`EOS`$) were gained DanielewiczSC298 . In particular, recent experimental and theoretical analyses were devoted to the study of the properties and effects of the symmetry term of the $`EOS`$ (asy-$`EOS`$) away from saturation conditions BaoAnBook01 ; baranRP .
Indeed, the two-component character of nuclear matter adds some special interest to the dynamics of heavy ion collisions at intermediate energies, between $`20`$ and $`100AMeV`$. In central collisions isospin distillation is an important effect in multifragmentation of charge asymmetric systems. Here phase separation is driven by isoscalar-like unstable fluctuations, i.e. local in phase variations of proton and neutron densities BaranPRL86 ; ColonnaPRL88 ; BaranNPA703 ; Jer . This leads to a more symmetric โliquidโ phase of fragments surrounded by a more neutron rich โgasโ relative to the original asymmetry of the system. Isoscaling phenomena, observed experimentally, provide indications for such a scenario XuPRL85 ; geraci .
In semi-peripheral collisions between nuclei with different $`N/Z`$ ratio, isospin dynamics will drive the system toward a uniform asymmetry distribution. The degree of equilibration, correlated to the interaction time, should provide some insights into transport properties of fermionic systems UehlingPR43 ; HellundPR56 , in particular give information on transport coefficients of asymmetric nuclear matter AndersonPRB35 ; ShiPRC68 .
The aim of this work is to investigate the isospin transfer through the neck region in semi-peripheral collisions of asymmetric nuclei at Fermi energies. Of particular interest is the role of the density dependence of the symmetry energy in this process. The isospin transfer was measured for collisions of different $`Sn`$ isotopes at MSU TsangNPA734 ; TsangPRL92 ; note1 and interpreted theoretically with the result that the asy-$`EOS`$ should be rather stiff. In these works the effect of pre-equilibrium emission, which changes the isospin content of the interacting system, was not analyzed explicitly. Here we will discuss this questions in detail, as well as the different transport processes affecting the final isospin content. Quantitatively, dynamical isospin effects can be properly understood only from microscopic calculations based on transport models. We will base our study on a stochastic BNV transport model (see refs. ColonnaNPA642 ; ColonnaRep for more details on the main ingredients of this approach).
## II ISOSPIN EQUILIBRATION PROCESS
We are focusing on the charge asymmetric collision $`{}_{}{}^{124}Sn+^{112}Sn`$, at $`50AMeV`$ bombarding energy, to which we refer as the mixed system, $`(M)`$. To investigate the density ($`\rho `$) dependence we consider here two representative parameterizations of the symmetry energy, $`E_{sym}(\rho ,I)/A=C_{sym}(\rho )I^2,I=(NZ)/A`$ : one showing a rapidly increasing behaviour with density, roughly proportional to $`\rho ^2`$ (asysuperstiff) and one where a kind of saturation is observed above normal density (soft, $`SKM^{}`$) (see Ref.baranRP ; BaranNPA703 for more detail).
The BNV simulations have been performed for semi-peripheral collisions at impact parameters $`b=6,8,9,10fm`$. In the last two cases the reaction has dominantly a binary character and the charge asymmetry of primary projectile (target)- like fragments, $`PLF`$ ( $`TLF`$), should provide the essential information about the isospin equilibration rate. At $`b=8fm`$ already about $`25\%`$ of the events are ternary. An IMF can be formed in the mid-velocity region by neck fragmentation BaranNPA730 . For more central events this mechanism becomes dominant: at $`b=6fm`$, one or two IMFโs are found in more than $`70\%`$ of events BaranNPA703 . The fragment formation in the neck region will influence the final isospin distribution of the PLF/TLF, and will render considerably more difficult the interpretation of the results. Thus in the following we select only binary events in our analysis. We define the average interaction time, $`t_c`$, as the time elapsed between the initial touching and the moment when $`PLF`$ and $`TLF`$ reseparate. From our simulations we obtain $`t_c140,120,100,80fm/c`$ for the impact parameters $`b=6,8,9,10fm`$, respectively. Four hundred events were calculated for each initial condition and for each asy-$`EOS`$.
Typical density contour plots, at $`b=8fm`$ and $`b=10fm`$, are shown in Figure 1. We note the dynamical evolution of the overlap region: driven by the fast leading motion of $`PL`$ and $`TL`$ prefragments, the formation of a lower density interface can be clearly observed after around $`40fm/c`$. An isospin migration, or transport, takes place during this transient configuration of two residues with densities close to the normal one, separated by a dilute neck region. In contrast, in deep-inelastic collisions at lower energies the isospin equilibration is driven by the N/Z difference between the interacting nuclei having a quite uniform density profile without a low density interface until separation FarineZPA339 .
We quantify the degree of equilibration by the isospin transport (imbalance) ratio RamiPRL84 , defined as:
$$R_i=\frac{2I_i^MI_i^HI_i^L}{I_i^HI_i^L}$$
(1)
Here $`i=P,T`$ stands for the projectile-like (target-like) fragment. The quantities $`I_i`$ refer to the isospin, $`I=(NZ)/A`$, or in general to any isospin dependent quantity, characterizing the fragments at separation time, for the mixed reaction $`(M,124+112)`$, the reactions between neutron rich ($`H,124+124`$), and between neutron poor nuclei ($`L,112+112`$), respectively. A value of $`R_i`$ approaching zero is an indication of a larger degree of equilibration. The extreme cases $`R_T=1`$ and $`R_P=1`$ correspond to the absence of any isospin transfer.
The results of the calculations are shown in Figure 2 for the dependence of $`R_{P/T}`$ on the interaction time $`t_c`$ for the asysoft (squares) and asysuperstiff (circles) $`EOS`$โs. The figure also shows the experimental values extracted in ref. TsangPRL92 ; note1 at semi-peripheral collisions at about $`b=8fm`$. We conclude, as in ref. TsangPRL92 , that an asystiff-like $`EOS`$ provides a better agreement with the experimental observations.
A significant difference between the two equations of state is evident for larger interaction times, i.e. smaller values of the impact parameter, $`b=6,8fm`$. The smaller values of isospin transport ratios for the asysoft $`EOS`$ point toward a faster equilibration rate. In refs. TsangPRL92 ; Chen04 an explanation was based on the observation that below normal density the asysoft $`EOS`$ has a larger value of the symmetry energy. Therefore an enhanced isospin equilibration will occur if the diffusion takes place at uniform lower density. We intend to show, that in fact the mechanism of charge equilibration is more complicated due to dynamical evolution of the reaction at these energies. Fast particle emission and density gradients will also play a role. In the next section we investigate more in detail the various influences on the isospin transfer process.
## III ISOSPIN SHARING AT FERMI ENERGIES
The isospin content of the two residues in a mixed collision system at separation time is determined by the interplay between the particle emission to the gas from each nucleus during the overlap and the transfer of nucleons through the neck. We thus write simple balance equations:
$`I_P={\displaystyle \frac{A_P^0}{A_P}}(I_P^0{\displaystyle \frac{A_{gP}}{A_P^0}}I_{gP}{\displaystyle \frac{A_{PT}}{A_P^0}}I_{PT}+{\displaystyle \frac{A_{TP}}{A_P^0}}I_{TP})`$ (2)
$`I_T={\displaystyle \frac{A_T^0}{A_T}}(I_T^0{\displaystyle \frac{A_{gT}}{A_T^0}}I_{gT}+{\displaystyle \frac{A_{PT}}{A_T^0}}I_{PT}{\displaystyle \frac{A_{TP}}{A_T^0}}I_{TP})`$ (3)
Here $`I_P`$ ($`I_T`$) and $`A_P`$ ($`A_T`$) are the $`PLF`$ ($`TLF`$) asymmetry and mass at separation, $`I_P^0`$ ($`I_T^0`$) and $`A_P^0`$ ($`A_T^0`$) the initial projectile (target) asymmetry and mass. Then $`I_{gP}`$ ($`I_{gT}`$), $`A_{gP}`$ ($`A_{gT}`$) are the asymmetries and masses of the projectile/target โgasโ, i.e. of the pre-equilibrium particles emitted by the projectile (target) during the interaction time. Finally $`I_{PT}`$ ($`I_{TP}`$), and $`A_{PT}`$ ($`A_{TP}`$) are the asymmetry and mass of all nucleons transferred from projectile (target) to target (projectile).
In Figure 3 we plot the time evolution of the quantities $`I_{gP}`$, $`I_{gT}`$ and $`A_{gP}`$, $`A_{gT}`$, as well as the values of $`I_{PT}`$ and $`I_{TP}`$ for the asysoft and the asysuperstiff $`EOS`$ for two impact parameters. We note that $`I_{gP}`$ is much larger than $`I_P^0`$. The same is true for the target but the difference is smaller. Thus the pre-equilibrium emission reduces the $`N/Z`$ difference between the two nuclei, competing with the transfer process. This is also clearly seen in Figure 4 where the time evolution of the projectile, target, and composite system apparent asymmetry, $`I_P^{app}`$, $`I_T^{app}`$ and $`I_C^{app}`$ is shown. The apparent asymmetry is calculated taking into account only the preequilibrium emission but not the isospin transfer inside the matter. A stronger variation is observed for the neutron rich projectile. The composite system asymmetry, $`I_C^{app}`$, is the limiting value toward which the two participants would evolve in complete isospin equilibrium. In contrast, we show in the figure the final asymmetries $`I_P`$ and $`I_T`$ including the transfer processes.
Comparing the results for the two $`EOS`$โs in Figs. 3 and 4 we see that a more neutron rich composition of pre-equilibrium emission is generated in the asysoft case because below normal density, from where most of the emitted nucleons originate, the neutrons (protons) are less (more) bound than for the asysuperstiff $`EOS`$. The differences between the two asy-$`EOS`$โs are strongly reduced at larger impact parameters, as seen in the results for $`b=10fm`$ in Figures 3 and 4, since the interaction times are much shorter.
We next show in Figure 5 the dependence of the quantities introduced above on interaction time, i.e. impact parameter. Here the projectile (target) and gas final asymmetries $`I_P`$ ,$`I_{gP}`$ ($`I_T`$, $`I_{gT}`$) are confronted with the transferred asymmetries $`I_{PT}`$ and $`I_{TP}`$. For the projectile, that is neutron rich, both pre-equilibrium emission and nucleon transfer drive the system toward a more symmetric configuration (left). The two processes have opposite effects and thus tend to compensate for the target (right). Therefore the projectile asymmetry has a more pronounced deviation from the corresponding initial value in comparison to the target.
We find a a clear dependence on the asy-$`EOS`$ of the isospin transferred between the two nuclei. In particular, for $`b=8fm(t_c=120fm/c)`$ we observe that
$`I_{PT}^{(asysuperstiff)}>I_{PT}^{(asysoft)}>I_P^0=0.192`$ (4)
$`I_{TP}^{(asysuperstiff)}>I_T^0=0.107>I_{TP}^{(asysoft)}`$ (5)
We will show in the next chapter that the origin of these inequalities lies in the existence of the low density interface and the density dependence of symmetry energy. The asysuperstiff $`EOS`$ favors the neutron migration toward the neck region from both participants. This explains why simultaneously $`I_{PT}^{(asysuperstiff)}>I_P^0`$ and $`I_{TP}^{(asysuperstiff)}>I_T^0`$. We will see that for the asysoft $`EOS`$ this effect is weakened.
## IV INTERPRETATION OF THE RESULTS
The above arguments can be made more explicit by considering that the proton and neutron migration is dictated by the spatial gradients of the corresponding chemical potentials $`\mu _{p/n}(\rho _p,\rho _n,T)`$ balian . The currents of the two species can be expressed as follows
$`j_n=ct\mu _n(\rho _p,\rho _n,T)=ct[\left({\displaystyle \frac{\mu _n}{\rho _n}}\right)_{\rho _p,T}\rho _n+`$
$`\left({\displaystyle \frac{\mu _n}{\rho _p}}\right)_{\rho _n,T}\rho _p]`$
$`j_p=ct\mu _p(\rho _p,\rho _n,T)=ct[\left({\displaystyle \frac{\mu _p}{\rho _n}}\right)_{\rho _p,T}\rho _n+`$
$`\left({\displaystyle \frac{\mu _p}{\rho _p}}\right)_{\rho _n,T}\rho _p],`$
where $`ct`$ is a constant. Rewriting these expressions in terms of Landau parameters, $`F_0^{qq^{}}`$ (baranRP and references therein), using:
$$N_q(T)\frac{\mu _q}{\rho _q^{}}=\delta _{qq^{}}+F_0^{qq^{}},q=n,pq^{}=n,p$$
(6)
(where $`N_q(T)`$ is the level density), we obtain:
$`j_n`$ $`=`$ $`{\displaystyle \frac{ct}{2}}([N_n^1(1+I)+f_n^\rho +If_n^I)]\rho +\rho [N_n^1+f_n^I]I)`$ (7)
$`=`$ $`D_n^\rho \rho D_n^II`$
$`j_p`$ $`=`$ $`{\displaystyle \frac{ct}{2}}([N_p^1(1I)+f_p^\rho If_p^I)]\rho \rho [N_p^1+f_p^I]I)`$ (8)
$`=`$ $`D_p^\rho \rho D_p^II,`$
where
$`f_q^\rho `$ $`=`$ $`N_q^1(F_0^{qq}+F_0^{qq^{}})={\displaystyle \frac{U_q}{\rho _q}}+{\displaystyle \frac{U_q}{\rho _q^{}}},`$ (9)
$`f_q^I`$ $`=`$ $`N_q^1(F_0^{qq}F_0^{qq^{}})={\displaystyle \frac{U_q}{\rho _q}}{\displaystyle \frac{U_q}{\rho _q^{}}};(qq^{}).`$ (10)
$`U_q`$ is the neutron or proton mean-field potential. The second lines in eqs. (7,8) define $`p`$ and $`n`$ drift and diffusion coefficients due to density and isospin gradients, $`D_q^\rho `$ and $`D_q^I`$, respectively. The terms $`N_q^1(1\pm I)`$ and $`f_q^\rho \pm If_q^I`$, that appear in the density drift coefficients $`D_q^\rho `$, can be expressed as
$$N_q^1(1\pm I)=2N^1(\rho ,T)\pm 4I\frac{C_{sym}^{kin}}{\rho }+O(I^2);$$
$$f_q^\rho \pm If_q^I=F(\rho )\pm 4I\frac{C_{sym}^{pot}}{\rho }+O(I^2),(+n,p),$$
(11)
where the function $`F(\rho )`$ depends only on the isoscalar part of the interaction. We have denoted by $`C_{sym}^{kin}`$ and $`C_{sym}^{pot}`$ the kinetic and the potential part of the symmetry energy coefficient $`C_{sym}`$, respectively. Combining eqs. (11), one can see that the isovector part of the nuclear interaction enters the coefficients $`D_q^\rho `$ through the derivative of the total symmetry energy $`C_{sym}`$. On the other hand the isospin diffusion coefficients $`D_q^I`$ are proportional to the quantity
$$\rho (N_q^1+f_q^I)=4[C_{sym}\pm I(\rho \frac{C_{sym}}{\rho }C_{sym})],$$
(12)
and depend, in leading order, on the value of the symmetry energy coeffcient $`C_{sym}`$.
Various particular situations can be derived from these relations. In symmetric nuclear matter $`D_n^I=D_p^I`$ and $`D_n^\rho =D_p^\rho `$. In the absence of density gradients the proton current will flow oppositely and equal in magnitude to the neutron current. On the other hand, for density gradients only, in asymmetric nuclear matter the proton and neutron currents may have the same direction but assume different values, inducing isospin gradients dist . Such a situation can be encountered in semi-pheripheral collisions between identical, charge asymmetric nuclei with the formation of a dilute intermediate region.
For the two asy-$`EOS`$โs we calculate the coefficients $`D_q^i`$, $`i=\rho ,I`$ and $`q=n,p`$. We plot the ratios $`R_q^i=D_q^{i,asysuperstiff}/D_q^{i,asysoft}`$ in Figure 6 as a function of the density for a fixed asymmetry $`I=0.2`$ . These values of the asymmetry and density are close to the physical conditions expected for the projectile or target region. The only negative coefficient is $`D_p^I`$. The isospin gradients, directed from the projectile to the neck and from the neck to the target, induce neutron and proton flows in opposite directions. However the ratios of the corresponding coefficients are quite close to unity for the two asy-$`EOS`$โs and therefore the effects are similar.
Since the density gradient is oriented from projectile and target residues to the neck, neutrons and protons migrate from higher toward lower density regions. Around and below saturation density $`R_n^\rho =D_n^{\rho ,asysuperstiff}/D_n^{\rho ,asysoft}>1`$ and $`R_p^\rho =D_p^{\rho ,asysuperstiff}/D_p^{\rho ,asysoft}<1`$. These inequalities suggest that more neutrons and less protons migrate from projectile toward neck in the case of asysuperstiff $`EOS`$ resulting in the formation of a more neutron rich intermediate region. This is due to the larger value of the derivative of the symmetry energy around normal density and is in agreement with the behavior observed in the numerical simulations, see eqs. (4,5).
We note that for asysuperstiff $`EOS`$, in spite of an enhanced isospin migration toward the neck at separation time, the projectile residue is more asymmetric in comparison to the asysoft case. One reason is that for the asysuperstiff $`EOS`$ during the pre-equilibrium emission not as many neutrons are removed as for the asysoft $`EOS`$. Also, as it was shown, for the asysuperstiff $`EOS`$ more asymmetric matter is also transferred from the target.
We now estimate the effect of isospin transport and pre-equilibrium emission on the transport ratios. Within the following approximations
$`{\displaystyle \frac{A_P^0}{A_P}}|_M{\displaystyle \frac{A_P^0}{A_P}}|_H{\displaystyle \frac{A_P^0}{A_P}}|_L;{\displaystyle \frac{A_{gP}}{A_P^0}}|_M{\displaystyle \frac{A_{gP}}{A_P^0}}|_H{\displaystyle \frac{A_{gP}}{A_P^0}}|_L;`$ (13)
$`A_{TP}A_{PT};I_{gP}^MI_{gP}^H;I_{gT}^MI_{gT}^L,`$ (14)
where the indices $`M,H,L`$ refer again to the mixed (124+112), the neutron rich $`(124+124)`$ and the neutron deficient $`(112+112)`$ systems, we arrive at a simplified expression for the isospin transport ratio for the projectile which shows explicitely the dependence on the isospin transport $`(I_{PT}I_{TP})`$ and on the pre-equilibrium emission $`(I_{gP}^HI_{gP}^L)`$:
$$R_P1\frac{2\frac{A_{PT}}{A_P^{0,H}}(I_{PT}I_{TP})}{I_P^{0,H}I_P^{0,L}\frac{A_{gP}}{A_P^{0,H}}(I_{gP}^HI_{gP}^L)}$$
(15)
With similar approximations the target isospin transport ratio can be expressed as:
$$R_T1+\frac{2\frac{A_{PT}}{A_T^{0,L}}(I_{PT}I_{TP})}{I_T^{0,H}I_T^{0,L}\frac{A_{gT}}{A_T^{0,L}}(I_{gT}^HI_{gT}^L)}$$
(16)
It is observed that the transport ratios depend on the difference $`I_{PT}I_{TP}`$ as expected. However, in contrast to what was assumed in refs. TsangPRL92 ; Chen04 , it is seen that they also depend on the pre-equilibrium emission which reduces the absolute value of the transport ratios. Both effects are smaller in the case of the asysuperstiff $`EOS`$, for which indeed a larger $`R`$ ratio is obtained.
## V CONCLUSIONS
In this work we have studied processes related to isospin equilibration in semipheripheral collisions at Fermi energies and their dependence on the symmetry term of the $`EOS`$. A special feature of these reactions is the development of a low density interface between the two residues. The neck region is controlling the proton and neutron currents and their direction. The presence of density gradients also affects the isospin exchange between projectile and target and we have shown that this is sensitive to the density dependence of the symmetry energy. The neutron to proton ratio emitted during the interaction stage is also influenced by the asy-$`EOS`$. The interplay between the two processes leads to a stronger equilibration for asy-soft $`EOS`$, as it is evidenced by the isospin transport (imbalance) ratio. Actually, in the asy-stiff case, a larger isospin transfer is observed, due to the presence of density gradients, directed from $`PLF`$ and $`TLF`$ towards the neck region. However, since we are studying binary processes, finally we observe a kind of compensation between the asymmetry of the matter transferred from projectile to target ($`I_{PT}`$) and from target to projectile($`I_{TP}`$). ยฟFrom this point of view, to put in better evidence effects due to the presence of density gradients, it would be more appropriate to study events where fragments originating from the neck region are also detected, possibly with their isospin content.
In the present study a rapidly increasing symmetry energy at subnormal densities (asysuperstiff) appears to be in better agreement with the existing data TsangPRL92 . More recent work also considered momentum dependent interactions, showing that the more repulsive character of the overall dynamics may reduce the symmetry energy stiffness required to reproduce the data Chen04 . The study of the interplay between the effects due to the isoscalar and isovector part of the $`EOS`$ on isospin transport observables deserves further attention.
In conclusion, charge equilibration measurements in semi-peripheral heavy ion collisions at Fermi energies provide new independent observables to study the poorly known density dependence of the symmetry term of the nuclear EOS. This is of interest for other properties of asymmetric matter, like neutron skin and isovector collective response in finite nuclei, and may also be important for neutron star crust structures baranRP .
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# Geometry of Optimal Control Problems and Hamiltonian Systems
### Preface
These notes are based on the mini-course given in June 2004 in Cetraro, Italy, in the frame of a C.I.M.E. school. Of course, they contain much more material that I could present in the 6 hours course. The goal was to give an idea of the general variational and dynamical nature of nice and powerful concepts and results mainly known in the narrow framework of Riemannian Geometry. This concerns Jacobi fields, Morseโs index formula, Levi Civita connection, Riemannian curvature and related topics.
I tried to make the presentation as light as possible: gave more details in smooth regular situations and referred to the literature in more complicated cases. There is an evidence that the results described in the notes and treated in technical papers we refer to are just parts of a united beautiful subject to be discovered on the crossroads of Differential Geometry, Dynamical Systems, and Optimal Control Theory. I will be happy if the course and the notes encourage some young ambitious researchers to take part in the discovery and exploration of this subject.
Acknowledgments. I would like to express my gratitude to Professor Gamkrelidze for his permanent interest to this topic and many inspiring discussions and to thank participants of the school for their surprising and encouraging will to work in the relaxing atmosphere of the Mediterranean resort.
## Part I Lagrange multipliersโ geometry
### 1 Smooth optimal control problems
In these lectures we discuss some geometric constructions and results emerged from the investigation of smooth optimal control problems. Weโll consider problems with integral costs and fixed endpoints. A standard formulation of such a problem is as follows: Minimize a functional
$$J_{t_0}^{t_1}(u())=\underset{t_0}{\overset{t_1}{}}\phi (q(t),u(t))๐t,$$
$`(1)`$
where
$$\dot{q}(t)=f(q(t),u(t)),u(t)U,t[t_0,t_1],$$
$`(2)`$
$`q(t_0)=q_0,q(t_1)=q_1`$. Here $`q(t)^n,U^k`$, a control function $`u()`$ is supposed to be measurable bounded while $`q()`$ is Lipschitzian; scalar function $`\phi `$ and vector function $`f`$ are smooth. A pair $`(u(),q())`$ is called an admissible pair if it satisfies differential equation (2) but may violate the boundary conditions.
We usually assume that Optimal Control Theory generalizes classical Calculus of Variations. Unfortunately, even the most classical geometric variational problem, the length minimization on a Riemannian manifold, cannot be presented in the just described way. First of all, even simplest manifolds, like spheres, are not domains in $`^n`$. This does not look as a serious difficulty: we slightly generalize original formulation of the optimal control problem assuming that $`q(t)`$ belongs to a smooth manifold $`M`$ instead of $`^n`$. Then $`\dot{q}(t)`$ is a tangent vector to $`M`$ i.e. $`\dot{q}(t)T_{q(t)}M`$ and we assume that $`f(q,u)T_qM,q,u.`$ Manifold $`M`$ is called the state space of the optimal control problem.
Now weโll try to give a natural formulation of the length minimization problem as an optimal control problem on a Riemannian manifold $`M`$. Riemannian structure on $`M`$ is (by definition) a family of Euclidean scalar products $`,_q`$ on $`T_qM,qM`$, smoothly depending on $`q`$. Let $`f_1(q),\mathrm{},f_n(q)`$ be an orthonormal basis of $`T_qM`$ for the Euclidean structure $`,_q`$ selected in such a way that $`f_i(q)`$ are smooth with respect to $`q`$. Then any Lipschitzian curve on $`M`$ satisfies a differential equation of the form:
$$\dot{q}=\underset{i=1}{\overset{n}{}}u_i(t)f_i(q),$$
$`(3)`$
where $`u_i()`$ are measurable bounded scalar functions. In other words, any Lipschitzian curve on $`M`$ is an admissible trajectory of the control system (3). The Riemannian length of the tangent vector $`\underset{i=1}{\overset{n}{}}u_if_i(q)`$ is $`\left(\underset{i=1}{\overset{n}{}}u_i^2\right)^{1/2}`$. Hence the length of a trajectory of system (3) defined on the segment $`[t_0,t_1]`$ is $`\mathrm{}(u())=_{t_0}^{t_1}\left(\underset{i=1}{\overset{n}{}}u_i^2(t)\right)^{1/2}๐t`$. Moreover, it is easy to derive from the CauchyโSchwarz inequality that the length minimization is equivalent to the minimization of the functional $`J_{t_0}^{t_1}(u())=_{t_0}^{t_1}\underset{i=1}{\overset{n}{}}u_i^2(t)dt`$. The length minimization problem is thus reduced to a specific optimal control problem on the manifold of the form (1), (2).
Unfortunately, what Iโve just written was wrong. It would be correct if we could select a smooth orthonormal frame $`f_i(q),qM,i=1,\mathrm{},n`$. Of course, we can always do it locally, in a coordinate neighborhood of $`M`$ but, in general, we cannot do it globally. We cannot do it even on the 2-dimensional sphere: you know very well that any continuous vector field on the 2-dimensional sphere vanishes somewhere. We thus need another more flexible formulation of a smooth optimal control problem.
Recall that a smooth locally trivial bundle over $`M`$ is a submersion $`\pi :VM`$, where all fibers $`V_q=\pi ^1(q)`$ are diffeomorphic to each other and, moreover, any $`qM`$ possesses a neighborhood $`O_q`$ and a diffeomorphism $`\mathrm{\Phi }_q:O_q\times V_q\pi ^1(O_q)`$ such that $`\mathrm{\Phi }_q(q^{},V_q)=V_q^{},q^{}O_q`$. In a less formal language one can say that a smooth locally trivial bundle is a smooth family of diffeomorphic manifolds $`V_q`$ (the fibers) parametrized by the points of the manifold $`M`$ (the base). Typical example is the tangent bundle $`TM=\underset{qM}{}T_qM`$ with the canonical projection $`\pi `$ sending $`T_qM`$ into $`q`$.
Definition. A smooth control system with the state space $`M`$ is a smooth mapping $`f:VTM`$, where $`V`$ is a locally trivial bundle over $`M`$ and $`f(V_q)T_qM`$ for any fiber $`V_q,qM`$. An admissible pair is a bounded<sup>1</sup><sup>1</sup>1the term โboundedโ means that the closure of the image of the mapping is compact measurable mapping $`v():[t_0,t_1]V`$ such that $`t\pi (v(t))=q(t)`$ is a Lipschitzian curve in $`M`$ and $`\dot{q}(t)=f(v(t))`$ for almost all $`t[t_0,t_1]`$. Integral cost is a functional $`J_{t_0}^{t_1}(v())=\underset{t_0}{\overset{t_1}{}}\phi (v(t))๐t`$, where $`\phi `$ is a smooth scalar function on $`V`$.
Remark. The above more narrow definition of an optimal control problem on $`M`$ was related to the case of a trivial bundle $`V=M\times U,V_q=\{q\}\times U`$. For the length minimization problem we have $`V=TM,f=\mathrm{Id},\phi (v)=v,v_q,vT_qM,qM`$.
Of course, any general smooth control system on the manifold $`M`$ is locally equivalent to a standard control system on $`^n`$. Indeed, any point $`qM`$ possesses a coordinate neighborhood $`O_q`$ diffeomorphic to $`^n`$ and a mapping $`\mathrm{\Phi }_q:O_q\times V_q\pi ^1(O_q)`$ trivializing the restriction of the bundle $`V`$ to $`O_q`$; moreover, the fiber $`V_q`$ can be embedded in $`^k`$ and thus serve as a set of control parameters $`U`$.
Yes, working locally we do not obtain new systems with respect of those in $`^n`$. Nevertheless, general intrinsic definition is very useful and instructive even for a purely local geometric analysis. Indeed, we do not need to fix specific coordinates on $`M`$ and a trivialization of $`V`$ when we study a control system defined in the intrinsic way. A change of coordinates in $`M`$ is actually a smooth transformation of the state space while a change of the trivialization results in the feedback transformation of the control system. This means that an intrinsically defined control system represents actually the whole class of systems that are equivalent with respect to smooth state and feedback transformations. All information on the system obtained in the intrinsic language is automatically invariant with respect to smooth state and feedback transformations. And this is what any geometric analysis intends to do: to study properties of the object under consideration preserved by the natural transformation group.
We denote by $`L_{\mathrm{}}([t_0,t_1];V)`$ the space of measurable bounded mappings from $`[t_0,t_1]`$ to $`V`$ equipped with the $`L_{\mathrm{}}`$-topology of the uniform convergence on a full measure subset of $`[t_0,t_1]`$. If $`V`$ would an Euclidean space, then $`L_{\mathrm{}}([t_0,t_1];V)`$ would have a structure of a Banach space. Since $`V`$ is only a smooth manifold, then $`L_{\mathrm{}}([t_0,t_1];V)`$ possesses a natural structure of a smooth Banach manifold modeled on the Banach space $`L_{\mathrm{}}([t_0,t_1];^{dimV})`$.
Assume that $`VM`$ is a locally trivial bundle with the $`n`$-dimensional base and $`m`$-dimensional fibers; then $`V`$ is an $`(n+m)`$-dimensional manifold.
###### Proposition I.1
Let $`f:VTM`$ be a smooth control system; then the space $`๐ฑ`$ of admissible pairs of this system is a smooth Banach submanifold of $`L_{\mathrm{}}([t_0,t_1];V)`$ modeled on $`^n\times L_{\mathrm{}}([t_0,t_1];^m)`$.
Proof. Let $`v()`$ be an admissible pair and $`q(t)=\pi (v(t)),t[t_0,t_1]`$. There exists a Lipschitzian with respect to $`t`$ family of local trivializations $`R_t:O_{q(t)}\times U\pi ^1(O_{q(t)})`$, where $`U`$ is diffeomorphic to the fibers $`V_q`$. The construction of such a family is a boring exercise which we omit.
Consider the system
$$\dot{q}=fR_t(q,u),uU.$$
$`(4)`$
Let $`v(t)=R_t(q(t),u(t))`$; then $`R_t,t_0tt_1,`$ induces a diffeomorphism of an $`L_{\mathrm{}}`$-neighborhood of $`(q(),u())`$ in the space of admissible pairs for (4) on a neighborhood of $`v()`$ in $`๐ฑ`$. Now fix $`\overline{t}[t_0,t_1]`$. For any $`\widehat{q}`$ close enough to $`q(\overline{t})`$ and any $`u^{}()`$ sufficiently close to $`u()`$ in the $`L_{\mathrm{}}`$-topology there exists a unique Lipschitzian path $`q^{}()`$ such that $`\dot{q}^{}(t)=fR_t(q^{}(t),u^{}(t))),t_0tt_1,q^{}(\overline{t})=\widehat{q}`$; moreover the mapping $`(\widehat{q},u^{}())q^{}()`$ is smooth. In other words, the Cartesian product of a neighborhood of $`q(\overline{t})`$ in $`M`$ and a neighborhood of $`u()`$ in $`L_{\mathrm{}}([t_0,t_1],U)`$ serves as a coordinate chart for a neighborhood of $`v()`$ in $`๐ฑ`$. This finishes the proof since $`M`$ is an $`n`$-dimensional manifold and $`L_{\mathrm{}}([t_0,t_1],U)`$ is a Banach manifold modeled on $`L_{\mathrm{}}([t_0,t_1],^m).\mathrm{}`$
An important role in our study will be played by the โevaluation mappingsโ $`F_t:v()q(t)=\pi (v(t))`$. It is easy to show that $`F_t`$ is a smooth mapping from $`๐ฑ`$ to $`M`$. Moreover, it follows from the proof of Proposition I.1 that $`F_t`$ is a submersion. Indeed, $`q(t)=F_t(v())`$ is, in fact a part of the coordinates of $`v()`$ built in the proof (the remaining part of the coordinates is the control $`u()`$.
### 2 Lagrange multipliers
Smooth optimal control problem is a special case of the general smooth conditional minimum problem on a Banach manifold $`๐ฒ`$. The general problem consists of the minimization of a smooth functional $`J:๐ฒ`$ on the level sets $`\mathrm{\Phi }^1(z)`$ of a smooth mapping $`\mathrm{\Phi }:๐ฒN`$, where $`N`$ is a finite-dimensional manifold. In the optimal control problem we have $`๐ฒ=๐ฑ,N=M\times M,\mathrm{\Phi }=(F_{t_0},F_{t_1})`$.
An efficient classical way to study the conditional minimum problem is the Lagrange multipliers rule. Let us give a coordinate free description of this rule. Consider the mapping
$$\overline{\mathrm{\Phi }}=(J,\mathrm{\Phi }):๐ฒ\times N,\overline{\mathrm{\Phi }}(w)=(J(w),\mathrm{\Phi }(w)),w๐ฒ.$$
It is easy to see that any point of the local conditional minimum or maximum (i.e. local minimum or maximum of $`J`$ on a level set of $`\mathrm{\Phi }`$) is a critical point of $`\overline{\mathrm{\Phi }}`$. I recall that $`w`$ is a critical point of $`\overline{\mathrm{\Phi }}`$ if the differential $`D_w\overline{\mathrm{\Phi }}:T_w๐ฒT_{\overline{\mathrm{\Phi }}(w)}\left(\times N\right)`$ is not a surjective mapping. Indeed, if $`D_w\overline{\mathrm{\Phi }}`$ would surjective then, according to the implicit function theorem, the image $`\overline{\mathrm{\Phi }}(O_w)`$ of an arbitrary neighborhood $`O_w`$ of $`w`$ would contain a neighborhood of $`\overline{\mathrm{\Phi }}(w)=(J(w),\mathrm{\Phi }(w))`$; in particular, this image would contain an interval $`((J(w)\epsilon ,J(w)+\epsilon ),\mathrm{\Phi }(w))`$ that contradicts the local conditional minimality or maximality of $`J(w)`$.
The linear mapping $`D_w\overline{\mathrm{\Phi }}`$ is not surjective if and only if there exists a nonzero linear form $`\overline{\mathrm{}}`$ on $`T_{\overline{\mathrm{\Phi }}(w)}\left(\times N\right)`$ which annihilates the image of $`D_w\overline{\mathrm{\Phi }}`$. In other words, $`\overline{\mathrm{}}D_w\overline{\mathrm{\Phi }}=0`$, where $`\overline{\mathrm{}}D_w\overline{\mathrm{\Phi }}:T_w๐ฒ`$ is the composition of $`D_w\overline{\mathrm{\Phi }}`$ and the linear form $`\overline{\mathrm{}}:T_{\overline{\mathrm{\Phi }}(w)}\left(\times N\right)`$.
We have $`T_{\overline{\mathrm{\Phi }}(w)}\left(\times N\right)=\times T_{\mathrm{\Phi }(w)}N`$. Linear forms on $`\left(\times N\right)`$ constitute the adjoint space $`\left(\times N\right)^{}=T_{\mathrm{\Phi }(w)}^{}N`$, where $`T_{\mathrm{\Phi }(w)}^{}N`$ is the adjoint space of $`T_{\mathrm{\Phi }(w)}M`$ (the cotangent space to $`M`$ at the point $`\mathrm{\Phi }(w)`$). Hence $`\mathrm{}=\nu \mathrm{}`$, where $`\nu ,\mathrm{}T_{\mathrm{\Phi }(w)}^{}N`$ and
$$\overline{\mathrm{}}D_w\overline{\mathrm{\Phi }}=(\nu \mathrm{})(d_wJ,D_w\mathrm{\Phi })=\nu d_wJ+\mathrm{}D_w\mathrm{\Phi }.$$
We obtain the equation
$$\nu d_wJ+\mathrm{}D_w\mathrm{\Phi }=0.$$
$`(5)`$
This is the Lagrange multipliers rule: if $`w`$ is a local conditional extremum, then there exists a nontrivial pair $`(\nu ,\mathrm{})`$ such that equation (5) is satisfied. The pair $`(\nu ,\mathrm{})`$ is never unique: indeed, if $`\alpha `$ is a nonzero real number, then the pair $`(\alpha \nu ,\alpha \mathrm{})`$ is also nontrivial and satisfies equation (5). So the pair is actually defined up to a scalar multiplier; it is natural to treat this pair as an element of the projective space $`\left(T_{\mathrm{\Phi }(w)}^{}N\right)`$ rather than an element of the linear space.
The pair $`(\nu ,\mathrm{})`$ which satisfies (5) is called the Lagrange multiplier associated to the critical point $`w`$. The Lagrange multiplier is called normal if $`\nu 0`$ and abnormal if $`\nu =0`$. In these lectures we consider only normal Lagrange multipliers, they belong to a distinguished coordinate chart of the projective space $`\left(T_{\mathrm{\Phi }(w)}^{}N\right)`$.
Any normal Lagrange multiplier has a unique representative of the form $`(1,\mathrm{})`$; then (5) is reduced to the equation
$$\mathrm{}D_w\mathrm{\Phi }=d_wJ.$$
$`(6)`$
The vector $`\mathrm{}T_{\mathrm{\Phi }(w)}^{}N`$ from equation (6) is also called a normal Lagrange multiplier (along with $`(1,\mathrm{})`$).
### 3 Extremals
Now we apply the Lagrange multipliers rule to the optimal control problem. We have $`\mathrm{\Phi }=(F_{t_0},F_{t_1}):๐ฑM\times M`$. Let an admissible pair $`v๐ฑ`$ be a critical point of the mapping $`(J_{t_0}^{t_1},\mathrm{\Phi })`$, the curve $`q(t)=\pi (v(t)),t_0tt_1`$ be the corresponding trajectory, and $`\mathrm{}T_{(q(t_0),q(t_1))}^{}(M\times M)`$ be a normal Lagrange multiplier associated to $`v()`$. Then
$$\mathrm{}D_v(F_{t_0},F_{t_1})=d_vJ_{t_0}^{t_1}.$$
$`(7)`$
We have $`T_{(q(t_0),q(t_1))}^{}(M\times M)=T_{q(t_0)}^{}M\times T_{q(t_1)}^{}M`$, hence $`\mathrm{}`$ can be presented in the form $`\mathrm{}=(\lambda _{t_0},\lambda _{t_1})`$, where $`\lambda _{t_i}T_{q(t_i)}^{}M,i=0,1`$. Equation (7) takes the form
$$\lambda _{t_1}D_vF_{t_1}\lambda _{t_0}D_vF_{t_0}=d_vJ_{t_0}^{t_1}.$$
$`(8)`$
Note that $`\lambda _{t_1}`$ in (8) is uniquely defined by $`\lambda _{t_0}`$ and $`v`$. Indeed, assume that $`\lambda _{t_1}^{}D_vF_{t_1}\lambda _{t_0}D_vF_{t_0}=d_vJ_{t_0}^{t_1}`$ for some $`\lambda _{t_1}^{}T_{q(t_1)}^{}M`$. Then $`(\lambda _{t_1}^{}\lambda _{t_1})D_vF_{t_1}=0`$. Recall that $`F_{t_1}`$ is a submersion, hence $`D_vF_{t_1}`$ is a surjective linear map and $`\lambda _{t_1}^{}\lambda _{t_1}=0`$.
###### Proposition I.2
Equality (8) implies that for any $`t[t_0,t_1]`$ there exists a unique $`\lambda _tT_{q(t)}^{}M`$ such that
$$\lambda _tD_vF_t\lambda _{t_0}D_vF_{t_0}=d_vJ_{t_0}^t$$
$`(9)`$
and $`\lambda _t`$ is Lipschitzian with respect to $`t`$.
Proof. The uniqueness of $`\lambda _t`$ follows from the fact that $`F_t`$ is a submersion as it was explained few lines above. Let us proof the existence. To do that we use the coordinatization of $`๐ฑ`$ introduced in the proof of Proposition I.1, in particular, the family of local trivializations $`R_t:O_{q(t)}\times U\pi ^1(O_{q(t)})`$. Assume that $`v(t)=R_t(q(t),u(t)),t_0tt_1`$, where $`v()`$ is the referenced admissible pair from (8).
Given $`\tau [t_0,t_1],\widehat{q}O_{q(\tau )}`$ let $`tQ_\tau ^t(\widehat{q})`$ be the solution of the differential equation $`\dot{q}=R_t(q,u(t))`$ which satisfies the condition $`Q_\tau ^\tau (\widehat{q})=\widehat{q}`$. In particular, $`Q_\tau ^t(q(\tau ))=q(t)`$. Then $`Q_\tau ^t`$ is a diffeomorphism of a neighborhood of $`q(\tau )`$ on a neighborhood of $`q(t)`$. We define a Banach submanifold $`๐ฑ_\tau `$ of the Banach manifold $`๐ฑ`$ in the following way:
$$๐ฑ_\tau =\{v^{}๐ฑ:\pi (v^{}(t))=Q_\tau ^t(\pi (v^{}(\tau ))),\tau tt_1\}.$$
It is easy to see that $`F_{t_1}|_{๐ฑ_\tau }=Q_\tau ^{t_1}F_\tau |_{๐ฑ_\tau }`$ and $`J_\tau ^{t_1}|_{๐ฑ_\tau }=a_\tau F_\tau `$, where $`a(\widehat{q})=\underset{\tau }{\overset{t}{}}\phi \left(\mathrm{\Phi }_t(Q_\tau ^t(\widehat{q}),u(t))\right)๐t`$. On the other hand, the set $`\{v^{}๐ฑ:v^{}|_{[t_0,\tau ]}๐ฑ_\tau |_{[t_0,\tau ]}\}`$ is a neighborhood of $`v`$ in $`๐ฑ`$. The restriction of (8) to $`๐ฑ_\tau `$ gives:
$$\lambda _{t_1}D_v\left(Q_\tau ^{t_1}F_\tau \right)\lambda _{t_0}D_vF_{t_0}=d_vJ_{t_0}^\tau +d_v\left(a_\tau F_\tau \right).$$
Now we apply the chain rule for the differentiation and obtain:
$$\lambda _\tau D_vF_\tau \lambda _{t_0}D_vF_{t_0}=d_vJ_{t_0}^\tau ,$$
where $`\lambda _\tau =\lambda _{t_1}D_{q(\tau )}Q_\tau ^{t_1}d_{q(\tau )}a_\tau `$. $`\mathrm{}`$
Definition. A Lipschitzian curve $`t\lambda _t,t_0tt_1,`$ is called a normal extremal of the given optimal control problem if there exists an admissible pair $`v๐ฑ`$ such that equality (9) holds. The projection $`q(t)=\pi (\lambda _t)`$ of a normal extremal is called a (normal) extremal path or a (normal) extremal trajectory.
According to Proposition I.2, normal Lagrange multipliers are just points of normal extremals. A good thing about normal extremals is that they satisfy a nice differential equation which links optimal control theory with a beautiful and powerful mathematics and, in many cases, allows to explicitly characterize all extremal paths.
### 4 Hamiltonian system
Here we derive equations which characterize normal extremals; we start from coordinate calculations. Given $`\tau [t_0,t_1]`$, fix a coordinate neighborhood $`๐ช`$ in $`M`$ centered at $`q(\tau )`$, and focus on the piece of the extremal path $`q()`$ which contains $`q()`$ and is completely contained in $`๐ช`$. Identity (9) can be rewritten in the form
$$\lambda _tD_vF_t\lambda _\tau D_vF_\tau =d_vJ_\tau ^t,$$
$`(10)`$
where $`q(t)`$ belongs to the piece of $`q()`$ under consideration. Fixing coordinates and a local trivialization of $`V`$ we (locally) identify our optimal control problem with a problem (1), (2) in $`^n`$. We have $`T^{}^n^n\times ^n=\{(p,q):p,q^n\}`$, where $`T_q^{}^n=^n\times \{q\}`$. Then $`\lambda _t=\{p(t),q(t)\}`$ and $`\lambda _tD_vF_t=p(t),D_vF_t=D_vp(t),F_t`$.
Admissible pairs of (2) are parametrized by $`\widehat{q}=F_\tau (v^{}),v^{}๐ฑ`$, and control functions $`u^{}()`$; the pairs have the form: $`v^{}=(u^{}(),q^{}(;\widehat{q},u^{}()))`$, where $`\frac{}{t}q^{}(t;\widehat{q},u^{}())=f(q^{}(t;\widehat{q},u^{}()),u^{}(t))`$ for all available $`t`$ and $`q^{}(\tau ;\widehat{q},u())=\widehat{q}`$. Then $`F_t(v^{})=q^{}(t;\widehat{q},u^{}())`$.
Now we differentiate identity (10) with respect to $`t`$: $`\frac{}{t}D_vp(t),F_t=\frac{}{t}d_vJ_\tau ^t`$ and change the order of the differentiation $`D_v\frac{}{t}p(t),F_t=d_v\frac{}{t}J_\tau ^t`$. We compute the derivatives with respect to $`t`$ at $`t=\tau `$:
$$\frac{}{t}p(t),F_t|_{t=\tau }=\dot{p}(\tau ),\widehat{q}+p(\tau ),f(\widehat{q},u^{}(\tau ),\frac{}{t}J_\tau ^t|_{t=\tau }=\phi (\widehat{q},u^{}(\tau )).$$
Now we have to differentiate with respect to $`v^{}()=(u^{}(),q^{}())`$. We however see that the quantities to differentiate depend only on the values of $`u^{}()`$ and $`q^{}()`$ at $`\tau `$, i.e. on the finite-dimensional vector $`(u^{}(\tau ),\widehat{q})`$. We derive:
$$\dot{p}(\tau )+\frac{}{q}p(\tau ),f(q(\tau ),u(\tau ))=\frac{\phi }{q}(q(t),u(t)),$$
$$\frac{}{u}p(\tau ),f(q(\tau ),u(\tau ))=\frac{\phi }{u}(q(\tau ),u(\tau )),$$
where $`v()=(q(),u())`$.
Of course, we can change $`\tau `$ and perform the differentiation at any available moment $`t`$. Finally, we obtain that (10) is equivalent to the identities
$$\dot{p}(t)+\frac{}{q}\left(p(t),f(q(t),u(t))\phi (q(t),u(t))\right)=0,$$
$$\frac{}{u}\left(p(t),f(q(t),u(t))\phi (q(t),u(t))\right)=0,$$
which can be completed by the equation $`\dot{q}=f(q(t),u(t))`$. We introduce a function $`h(p,q,u)=p,f(q,u)\phi (q,u)`$ which is called the Hamiltonian of the optimal control problem (1), (2). This function permits us to present the obtained relations in a nice Hamiltonian form:
$$\{\begin{array}{cc}\hfill \dot{p}& =\frac{h}{q}(p,q,u)\hfill \\ \hfill \dot{q}& =\frac{h}{p}(p,q,u)\hfill \end{array},\frac{h}{u}(p,q,u)=0.$$
$`(11)`$
A more important fact is that system (11) has an intrinsic coordinate free interpretation. Recall that in the triple $`(p,q,u)`$ neither $`p`$ nor $`u`$ has an intrinsic meaning; the pair $`(p,q)`$ represents $`\lambda T^{}M`$ while the pair $`(q,u)`$ represents $`vV`$. First we consider an intermediate case $`V=M\times U`$ (when $`u`$ is separated from $`q`$ but coordinates in $`M`$ are not fixed) and then turn to the completely intrinsic setting.
If $`V=M\times U`$, then $`f:M\times UTM`$ and $`f(q,u)T_qM`$. The Hamiltonian of the optimal control problem is a function $`h:T^{}M\times U`$ defined by the formula $`h(\lambda ,u)=\lambda (f(q,u))\phi (q,u)`$, $`\lambda T_q^{}M,qM,uU`$. For any $`uU`$ we obtain a function $`h_u\stackrel{def}{=}h(,u)`$ on $`T^{}M`$. The cotangent bundle $`T^{}M`$ possesses a canonical symplectic structure which provides a standard way to associate a Hamiltonian vector field to any smooth function on $`T^{}M`$. Weโll recall this procedure.
Let $`\pi :T^{}MM`$ be the projection, $`\pi (T_q^{}M)=\{q\}`$. The Liouville (or tautological) differential 1-form $`\varsigma `$ on $`T^{}M`$ is defined as follows. Let $`\varsigma _\lambda :T_\lambda (T^{}M)`$ be the value of $`\varsigma `$ at $`\lambda T^{}M`$, then $`\varsigma _\lambda =\lambda \pi _{}`$, the composition of $`\pi _{}:T_\lambda (T^{}M)T_{\pi (\lambda )}M`$ and the cotangent vector $`\lambda :T_{\pi (\lambda )}M`$. The coordinate presentation of the Liouville form is: $`\varsigma _{(p,q)}=p,dq=\underset{i=1}{\overset{n}{}}p^idq^i`$, where $`p=(p^1,\mathrm{},p^n)`$, $`q=(q^1,\mathrm{},q^n)`$. The canonical symplectic structure on $`T^{}M`$ is the differential 2-form $`\sigma =d\varsigma `$; its coordinate representation is: $`\sigma =\underset{i=1}{\overset{n}{}}dp^idq^i`$. The Hamiltonian vector field associated to a smooth function $`a:T^{}M`$ is a unique vector field $`\stackrel{}{a}`$ on $`T^{}M`$ which satisfies the equation $`\sigma (,\stackrel{}{a})=da`$. The coordinate representation of this field is: $`\stackrel{}{a}=\underset{i=1}{\overset{n}{}}\left(\frac{a}{p_i}\frac{}{q_i}\frac{a}{q_i}\frac{}{p_i}\right)`$. Equations (11) can be rewritten in the form:
$$\dot{\lambda }=\stackrel{}{h}_u(\lambda ),\frac{h}{u}(\lambda ,u)=0.$$
$`(12)`$
Now let $`V`$ be an arbitrary locally trivial bundle over $`M`$. Consider the Cartesian product of two bundles:
$$T^{}M\times _MV=\{(\lambda ,v):vV_q,\lambda T_q^{}M,qM\}$$
that is a bundle over $`M`$ whose fibers are Cartesian products of the correspondent fibers of $`V`$ and $`T^{}M`$. Hamiltonian of the optimal control problem takes the form $`h(\lambda ,v)=\lambda (f(v))\phi (v)`$; this is a well-defined smooth function on $`T^{}M\times _MU`$. Let $`๐ญ:T^{}M\times _MVT^{}M`$ be the projection on the first factor, $`๐ญ:(\lambda ,v)\lambda `$. Equations (11) (or (12)) can be rewritten in the completely intrinsic form as follows: $`(๐ญ^{}\sigma )_v(,\dot{\lambda })=dh`$. One may check this fact in any coordinates; we leave this simple calculation to the reader.
Of course, by fixing a local trivialization of $`V`$, we turn the last relation back into a more convinient to study equation (12). A domain $`๐`$ in $`T^{}M`$ is called regular for the Hamiltonian $`h`$ if for any $`\lambda ๐`$ there exists a unique solution $`u=\overline{u}(\lambda )`$ of the equation $`\frac{h}{u}(\lambda ,u)=0`$, where $`\overline{u}(\lambda )`$ is smooth with respect to $`\lambda `$. In particular, if $`U`$ is an affine space and the functions $`uh(\lambda ,u)`$ are strongly concave (convex) and bounded from above (below) for $`\lambda ๐`$, then $`๐`$ is regular and $`\overline{u}(\lambda )`$ is defined by the relation
$$h(\lambda ,\overline{u}(\lambda ))=\underset{uU}{\mathrm{max}}h(\lambda ,u)\left(h(\lambda ,\overline{u}(\lambda ))=\underset{uU}{\mathrm{min}}h(\lambda ,u)\right).$$
In the regular domain, we set $`H(\lambda )=h(\lambda ,\overline{u}(\lambda ))`$, where $`\frac{h}{u}(\lambda ,\overline{u}(\lambda ))=0`$. It is easy to see that equations (12) are equivalent to one Hamiltonian system $`\dot{\lambda }=\stackrel{}{H}(\lambda )`$. Indeed, the equality $`d_{(\lambda ,\overline{u}(\lambda ))}h=d_\lambda h_{\overline{u}(\lambda )}+\frac{h_{\overline{u}(\lambda )}}{u}du=d_\lambda h_{\overline{u}(\lambda )}`$ immediately implies that $`\stackrel{}{H}(\lambda )=\stackrel{}{h}_{\overline{u}(\lambda )}(\lambda )`$.
### 5 Second order information
We come back to the general setting of Section 2 and try to go beyond the Lagrange multipliers rule. Take a pair $`(\mathrm{},w)`$ which satisfies equation (6). We call such pairs (normal) Lagrangian points. Let $`\mathrm{\Phi }(w)=z`$. If $`w`$ is a regular point of $`\mathrm{\Phi }`$, then $`\mathrm{\Phi }^1(z)O_w`$ is a smooth codimension $`dimN`$ submanifold of $`๐ฒ`$, for some neighborhood $`O_w`$ of $`w`$. In this case $`w`$ is a critical point of $`J|_{\mathrm{\Phi }^1(z)O_w}`$. We are going to compute the Hessian of $`J|_{\mathrm{\Phi }^1(z)}`$ at $`w`$ without resolving the constraints $`\mathrm{\Phi }(w)=z`$. The formula we obtain makes sense without the regularity assumptions as well.
Let $`s\gamma (s)`$ be a smooth curve in $`\mathrm{\Phi }^1(z)`$ such that $`\gamma (0)=w`$. Differentiation of the identity $`\mathrm{\Phi }(\gamma (s))=z`$ gives:
$$D_w\mathrm{\Phi }\dot{\gamma }=0,D_w^2\mathrm{\Phi }(\dot{\gamma },\dot{\gamma })+D_w\mathrm{\Phi }\ddot{\gamma }=0,$$
where $`\dot{\gamma }`$ and $`\ddot{\gamma }`$ are the first and the second derivatives of $`\gamma `$ at $`s=0`$. We also have:
$$\frac{d^2}{ds^2}J(\gamma (s))|_{s=0}=D_w^2J(\dot{\gamma },\dot{\gamma })+D_wJ\ddot{\gamma }\stackrel{\mathrm{eq}.(6)}{=}$$
$$D_w^2J(\dot{\gamma },\dot{\gamma })+\mathrm{}D_w\mathrm{\Phi }\ddot{\gamma }=D_w^2J(\dot{\gamma },\dot{\gamma })\mathrm{}D_w^2\mathrm{\Phi }(\dot{\gamma },\dot{\gamma }).$$
Finally,
$$\mathrm{Hess}_w(J|_{\mathrm{\Phi }^1(z)})=(D_w^2J\mathrm{}D_w^2\mathrm{\Phi })|_{\mathrm{ker}D_w\mathrm{\Phi }}.$$
$`(13)`$
###### Proposition I.3
If quadratic form (13) is positive (negative) definite, then $`w`$ is a strict local minimizer (maximizer) of $`J|_{\mathrm{\Phi }^1(z)}`$.
If $`w`$ is a regular point of $`\mathrm{\Phi }`$, then the proposition is obvious but one can check that it remains valid without the regularity assumption. On the other hand, without the regularity assumption, local minimality does not imply nonnegativity of form (13). What local minimality (maximality) certainly implies is nonnegativity (nonpositivity) of form (13) on a finite codimension subspace of $`\mathrm{ker}D_w\mathrm{\Phi }`$ (see \[7, Ch. 20\] and references there).
Definition. A Lagrangian point $`(\mathrm{},w)`$ is called sharp if quadratic form (13) is nonnegative or nonpositive on a finite codimension subspace of $`\mathrm{ker}D_w\mathrm{\Phi }`$.
Only sharp Lagrangian points are counted in the conditional extremal problems under consideration. Let $`Q`$ be a real quadratic form defined on a linear space $`E`$. Recall that the negative inertia index (or the Morse index) $`\mathrm{ind}Q`$ is the maximal possible dimension of a subspace in $`E`$ such that the restriction of $`Q`$ to the subspace is a negative form. The positive inertia index of $`Q`$ is the Morse index of $`Q`$. Each of these indices is a nonnegative integer or $`+\mathrm{}`$. A Lagrangian point $`(\mathrm{},w)`$ is sharp if the negative or positive inertia index of form (13) is finite.
In the optimal control problems, $`๐ฒ`$ is a huge infinite dimensional manifold while $`N`$ usually has a modest dimension. It is much simpler to characterize Lagrange multipliers in $`T^{}N`$ (see the previous section) than to work directly with $`J|_{\mathrm{\Phi }^1(z)}`$. Fortunately, the information on the sign and, more generally, on the inertia indices of the infinite dimensional quadratic form (13) can also be extracted from the Lagrange multipliers or, more precisely, from the so called $``$-derivative that can be treated as a dual to the form (13) object.
$``$-derivative concerns the linearization of equation (6) at a given Lagrangian point. In order to linearize the equation we have to present its left- and right-hand sides as smooth mappings of some manifolds. No problem with the right-hand side: $`wd_wJ`$ is a smooth mapping from $`๐ฒ`$ to $`T^{}๐ฒ`$. The variables $`(\mathrm{},w)`$ of the left-hand side live in the manifold
$$\mathrm{\Phi }^{}T^{}N=\{(\mathrm{},w):\mathrm{}T_{\mathrm{\Phi }(w)}^{},w๐ฒ\}T^{}N\times ๐ฒ.$$
Note that $`\mathrm{\Phi }^{}T^{}N`$ is a locally trivial bundle over $`๐ฒ`$ with the projector $`\pi :(\mathrm{},w)w`$; this is nothing else but the induced bundle from $`T^{}N`$ by the mapping $`\mathrm{\Phi }`$. We treat equation (6) as the equality of values of two mappings from $`\mathrm{\Phi }^{}T^{}N`$ to $`T^{}๐ฒ`$. Let us rewrite this equation in local coordinates.
So let $`N=^m`$ and $`๐ฒ`$ be a Banach space. Then $`T^{}N=^m\times ^m`$ (where $`T_zN=^m\times \{z\}`$), $`T^{}๐ฒ=๐ฒ^{}\times ๐ฒ`$, $`\mathrm{\Phi }^{}T^{}N=^m\times ^m\times ๐ฒ`$. Surely, $`^m^m`$ but in the forthcoming calculations it is convenient to treat the first factor in the product $`^m\times ^m`$ as the space of linear forms on the second factor. We have: $`\mathrm{}=(\zeta ,z)^m\times ^m`$ and equation (6) takes the form
$$\zeta \frac{d\mathrm{\Phi }}{dw}=\frac{dJ}{dw},\mathrm{\Phi }(w)=z.$$
$`(14)`$
Linearization of system (14) at the point $`(\zeta ,z,w)`$ reads:
$$\zeta ^{}\frac{d\mathrm{\Phi }}{dw}+\zeta \frac{d^2\mathrm{\Phi }}{dw^2}(w^{},)=\frac{d^2J}{dw^2}(w^{},),\frac{d\mathrm{\Phi }}{dw}w^{}=z^{}.$$
$`(15)`$
We set
$$_{(\mathrm{},w)}^0(\overline{\mathrm{\Phi }})=\{\mathrm{}^{}=(\zeta ^{},z^{})T_{\mathrm{}}(T^{}N):w^{}๐ฒ\mathrm{s}.\mathrm{t}.(\zeta ^{},z^{},w^{})\mathrm{satisfies}(15)\}.$$
Note that subspace $`_{(\mathrm{},w)}^0(\overline{\mathrm{\Phi }})T_{\mathrm{}}(T^{}N)`$ does not depend on the choice of local coordinates. Indeed, to construct this subspace we take all $`(\mathrm{}^{},w^{})T_{(\mathrm{},w)}(\mathrm{\Phi }^{}T^{}N)`$ which satisfy the linearized equation (6) and then apply the projection $`(\mathrm{}^{},w^{})\mathrm{}^{}`$.
Recall that $`T_{\mathrm{}}(T^{}N)`$ is a symplectic space endowed with the canonical symplectic form $`\sigma _{\mathrm{}}`$ (cf. Sec. 4). A subspace $`ST_{\mathrm{}}(T^{}N)`$ is isotropic if $`\sigma _{\mathrm{}}|_S=0`$. Isotropic subspaces of maximal possible dimension $`m=\frac{1}{2}dimT_{\mathrm{}}(T^{}N)`$ are called Lagrangian subspaces.
###### Proposition I.4
$`_{(\mathrm{},w)}^0(\overline{\mathrm{\Phi }})`$ is an isotropic subspace of $`T_{\mathrm{}}(T^{}N)`$. If $`dim๐ฒ<\mathrm{}`$, then $`_{(\mathrm{},w)}^0(\overline{\mathrm{\Phi }})`$ is a Lagrangian subspace.
Proof. First weโll prove the isotropy of $`_{(\mathrm{},w)}^0(\overline{\mathrm{\Phi }})`$. Let $`(\zeta ^{},z^{}),(\zeta ^{\prime \prime },z^{\prime \prime })T_{\mathrm{}}(T^{}N)`$. We have $`\sigma _{\mathrm{}}((\zeta ^{},z^{}),(\zeta ^{\prime \prime },z^{\prime \prime }))=\zeta ^{}z^{\prime \prime }\zeta ^{\prime \prime }z^{}`$; here the symbol $`\zeta z`$ denotes the result of the application of the linear form $`\zeta ^m`$ to the vector $`z^n`$ or, in the matrix terminology, the product of the row $`\zeta `$ and the column $`z`$. Assume that $`(\zeta ^{},z^{},w^{})`$ and $`(\zeta ^{\prime \prime },z^{\prime \prime },w^{\prime \prime })`$ satisfy equations (15); then
$$\zeta ^{}z^{\prime \prime }=\zeta ^{}\frac{d\mathrm{\Phi }}{dw}w^{\prime \prime }=\frac{d^2J}{dw^2}(w^{},w^{\prime \prime })\zeta \frac{d^2\mathrm{\Phi }}{dw^2}(w^{},w^{\prime \prime }).$$
$`(16)`$
The right-hand side of (16) is symmetric with respect to $`w^{}`$ and $`w^{\prime \prime }`$ due to the symmetry of second derivatives. Hence $`\zeta ^{}z^{\prime \prime }=\zeta ^{\prime \prime }z^{}`$. In other words, $`\sigma _{\mathrm{}}((\zeta ^{},z^{}),(\zeta ^{\prime \prime },z^{\prime \prime }))=0`$. So $`_{(\mathrm{},w)}^0(\overline{\mathrm{\Phi }})`$ is isotropic and, in particular, $`dim\left(_{(\mathrm{},w)}^0(\overline{\mathrm{\Phi }})\right)m`$.
Now show that the last inequality becomes the equality as soon as $`๐ฒ`$ is finite dimensional. Set $`Q=\frac{d^2J}{dw^2}\zeta \frac{d^2\mathrm{\Phi }}{dw^2}`$ and consider the diagram:
$$\zeta ^{}\frac{d\mathrm{\Phi }}{dw}Q(w^{},)\stackrel{left}{}(\zeta ^{},w^{})\stackrel{right}{}(\zeta ^{},\frac{d\mathrm{\Phi }}{dw}w^{}).$$
Then $`_{(\mathrm{},w)}^0(\overline{\mathrm{\Phi }})=right(\mathrm{ker}(left))`$. Passing to a factor space if necessary we may assume that $`\mathrm{ker}(left)\mathrm{ker}(right)=0`$; this means that:
$$\frac{d\mathrm{\Phi }}{dw}w^{}\&Q(w^{},)=0w^{}=0.$$
$`(17)`$
Under this assumption, $`dim_{(\mathrm{},w)}^0(\overline{\mathrm{\Phi }})=dim\mathrm{ker}(left)`$. On the other hand, relations (17) imply that the mapping $`left:^m\times ๐ฒ๐ฒ^{}`$ is surjective. Indeed, if, on the contrary, the map $`left`$ is not surjective then there exists a nonzero vector $`v(๐ฒ^{})^{}=๐ฒ`$ which annihilates the image of $`left`$; in other words, $`\zeta ^{}\frac{d\mathrm{\Phi }}{dw}vQ(w^{},v)=0,\zeta ^{},w^{}`$. Hence $`\frac{d\mathrm{\Phi }}{dw}v=0\&Q(v,)=0`$ that contradicts (17). It follows that $`dim_{(\mathrm{},w)}^0(\overline{\mathrm{\Phi }})=dim(^m\times ๐ฒ)dim๐ฒ^{}=m.\mathrm{}`$
For infinite dimensional $`๐ฒ`$, the space $`_{(\mathrm{},w)}^0(\overline{\mathrm{\Phi }})`$ may have dimension smaller than $`m`$ due to an ill-posedness of equations (15); to guarantee dimension $`m`$ one needs certain coercivity of the form $`\zeta \frac{d^2\mathrm{\Phi }}{dw^2}`$. I am not going to discuss here what kind of coercivity is sufficient, it can be easily reconstructed from the proof of Proposition I.4 (see also ). Anyway, independently on any coercivity one can take a finite dimensional approximation of the original problem and obtain a Lagrangian subspace $`_{(\mathrm{},w)}^0(\overline{\mathrm{\Phi }})`$ guaranteed by Proposition I.4. What happens with these subspaces when the approximation becomes better and better, do they have a well-defined limit (which would be unavoidably Lagrangian)? A remarkable fact is that such a limit does exist for any sharp Lagrangian point. It contains $`_{(\mathrm{},w)}^0(\overline{\mathrm{\Phi }})`$ and is called the $``$-derivative of $`\overline{\mathrm{\Phi }}`$ at $`(\mathrm{},w)`$. To formulate this result we need some basic terminology from set theoretic topology.
A partially ordered set $`(๐,)`$ is a directed set if $`\alpha _1,\alpha _2๐`$ $`\beta ๐`$ such that $`\alpha _1\beta `$ and $`\alpha _2\beta `$. A family $`\{x_\alpha \}_{\alpha ๐}`$ of points of a topological space $`๐ณ`$ indexed by the elements of $`๐`$ is a generalized sequence in $`๐ณ`$. A point $`x๐ณ`$ is the limit of the generalized sequence $`\{x_\alpha \}_{\alpha ๐}`$ if for any neighborhood $`๐ช_x`$ of $`x`$ in $`๐ณ`$ $`\alpha ๐`$ such that $`x_\beta ๐ช_x,\beta \alpha `$; in this case we write $`x=\underset{๐}{lim}x_\alpha `$.
Let $`๐ด`$ be a finite dimensional submanifold of $`๐ฒ`$ and $`w๐ด`$. If $`(\mathrm{},w)`$ is a Lagrangian point for $`\overline{\mathrm{\Phi }}=(J,\mathrm{\Phi })`$, then it is a Lagrangian point for $`\overline{\mathrm{\Phi }}|_๐ด`$. A straightforward calculation shows that the Lagrangian subspace $`_{(\mathrm{},w)}^0(\overline{\mathrm{\Phi }}|_๐ด)`$ depends on the tangent space $`W=T_w๐ด`$ rather than on $`๐ด`$, i.e. $`_{(\mathrm{},w)}^0(\overline{\mathrm{\Phi }}|_๐ด)=_{(\mathrm{},w)}^0(\overline{\mathrm{\Phi }}|_๐ด^{})`$ as soon as $`T_w๐ด=T_w๐ด^{}=W`$. We denote $`\mathrm{\Lambda }_W=_{(\mathrm{},w)}^0(\overline{\mathrm{\Phi }}|_๐ด)`$. Recall that $`\mathrm{\Lambda }_W`$ is an $`m`$-dimensional subspace of the $`2m`$-dimensional space $`T_{\mathrm{}}(T^{}N)`$, i.e. $`\mathrm{\Lambda }_W`$ is a point of the Grassmann manifold of all $`m`$-dimensional subspaces in $`T_{\mathrm{}}(T^{}N)`$.
Finally, we denote by $`๐`$ the set of all finite dimensional subspaces of $`T_w๐ฒ`$ partially ordered by the inclusion โ$``$โ. Obviously, $`(๐,)`$ is a directed set and $`\{\mathrm{\Lambda }_W\}_{W๐}`$ is a generalized sequence indexed by the elements of this directed set. It is easy to check that there exists $`W_0๐`$ such that $`\mathrm{\Lambda }_W_{(\mathrm{},w)}^0(\overline{\mathrm{\Phi }}),WW_0`$. In particular, if $`_{(\mathrm{},w)}^0(\overline{\mathrm{\Phi }})`$ is $`m`$-dimensional, then $`\mathrm{\Lambda }_{W_0}=_{(\mathrm{},w)}^0(\overline{\mathrm{\Phi }}),WW_0`$, the sequence $`\mathrm{\Lambda }_W`$ is stabilizing and $`_{(\mathrm{},w)}^0(\overline{\mathrm{\Phi }})=\underset{๐}{lim}\mathrm{\Lambda }_W`$. In general, the sequence $`\mathrm{\Lambda }_W`$ is not stabilizing, nevertheless the following important result is valid.
###### Theorem I.1
If $`(\mathrm{},w)`$ is a sharp Lagrangian point, then there exists $`_{(\mathrm{},w)}(\overline{\mathrm{\Phi }})=\underset{๐}{lim}\mathrm{\Lambda }_W`$.
We omit the proof of the theorem, you can find this proof in paper with some other results which allow to efficiently compute $`\underset{๐}{lim}\mathrm{\Lambda }_W`$. Lagrangian subspace $`_{(\mathrm{},w)}(\overline{\mathrm{\Phi }})=\underset{๐}{lim}\mathrm{\Lambda }_W`$ is called the $``$-derivative of $`\overline{\mathrm{\Phi }}=(J,\mathrm{\Phi })`$ at the Lagrangian point $`(\mathrm{},w)`$.
Obviously, $`_{(\mathrm{},w)}(\overline{\mathrm{\Phi }})_{(\mathrm{},w)}^0(\overline{\mathrm{\Phi }})`$. One should think on $`_{(\mathrm{},w)}(\overline{\mathrm{\Phi }})`$ as on a completion of $`_{(\mathrm{},w)}^0(\overline{\mathrm{\Phi }})`$ by means of a kind of weak solutions to system (15) which could be missed due to the ill-posedness of the system.
Now we should explain the connection between $`_{(\mathrm{},w)}(\overline{\mathrm{\Phi }})`$ and $`\mathrm{Hess}_w(J|_{\mathrm{\Phi }^1(z)})`$. We start from the following simple observation:
###### Lemma I.1
Assume that $`dim๐ฒ<\mathrm{}`$, $`w`$ is a regular point of $`\mathrm{\Phi }`$ and $`\mathrm{ker}D_w\mathrm{\Phi }\mathrm{ker}(D_w^2J\mathrm{}D_w^2\mathrm{\Phi })=0`$. Then
$$\mathrm{ker}\mathrm{Hess}_w(J|_{\mathrm{\Phi }^1(z)})=0_{(\mathrm{},w)}(\overline{\mathrm{\Phi }})T_{\mathrm{}}(T_z^{}N)=0,$$
i.e. quadratic form $`\mathrm{Hess}_w(J|_{\mathrm{\Phi }^1(z)})`$ is nondegenerate if and only if the subspace $`_{(\mathrm{},w)}(\overline{\mathrm{\Phi }})`$ is transversal to the fiber $`T_z^{}N`$.
Proof. We make computations in coordinates. First, $`T_{\mathrm{}}(T_z^{}N)=\{(\zeta ^{},0):\zeta ^{}^n\}`$; then, according to equations (15), $`(\zeta ^{},0)_{(\mathrm{},w)}(\overline{\mathrm{\Phi }})`$ if and only if there exists $`w๐ฒ`$ such that
$$\frac{d\mathrm{\Phi }}{dw}w^{}=0,\frac{d^2J}{dw^2}(w^{},)\mathrm{}\frac{d^2\mathrm{\Phi }}{dw^2}(w^{},)=\zeta ^{}\frac{d\mathrm{\Phi }}{dw}.$$
$`(18)`$
Regularity of $`w`$ implies that $`\zeta ^{}\frac{d\mathrm{\Phi }}{dw}0`$ and hence $`w^{}0`$ as soon as $`\zeta ^{}0`$. Equalities (18) imply: $`\frac{d^2J}{dw^2}(w^{},v)\mathrm{}\frac{d^2\mathrm{\Phi }}{dw^2}(w^{},v)=0,v\mathrm{ker}\frac{d\mathrm{\Phi }}{dw}`$, i.e. $`w^{}\mathrm{ker}\mathrm{Hess}_w(J|_{\mathrm{\Phi }^1(z)})`$. Moreover, our implications are invertible: we could start from a nonzero vector $`w^{}\mathrm{ker}\mathrm{Hess}_w(J|_{\mathrm{\Phi }^1(z)})`$ and arrive to a nonzero vector $`(\zeta ^{},0)_{(\mathrm{},w)}(\overline{\mathrm{\Phi }}).\mathrm{}`$
Remark. Condition $`\mathrm{ker}D_w\mathrm{\Phi }\mathrm{ker}(D_w^2J\mathrm{}D_w^2\mathrm{\Phi })=0`$ from Lemma I.1 is not heavy. Indeed, a pair $`(J,\mathrm{\Phi })`$ satisfies this condition at all its Lagrangian points if and only if 0 is a regular value of the mapping $`(\zeta ,w)\zeta \frac{d\mathrm{\Phi }}{dw}\frac{dJ}{dw}`$. Standard Transversality Theorem implies that this is true for generic pair $`(J,\mathrm{\Phi })`$.
### 6 Maslov index
Lemma I.1 is a starting point for a far going theory which allows to effectively compute the Morse index of the Hessians in terms of the $``$-derivatives.
How to do it? Normally, extremal problems depend on some parameters. Actually, $`zN`$ is such a parameter and there could be other ones, which we do not explicitly add to the constraints. In the optimal control problems a natural parameter is the time interval $`t_1t_0`$. Anyway, assume that we have a continuous family of the problems and their sharp Lagrangian points: $`\mathrm{}_\tau D_{w_\tau }\mathrm{\Phi }_\tau =d_{w_\tau }J_\tau ,\tau _0\tau \tau _1`$; let $`\mathrm{\Lambda }(\tau )=_{(\mathrm{}_\tau ,w_\tau )}(\overline{\mathrm{\Phi }}_\tau )`$. Our goal is to compute the difference $`\mathrm{ind}\mathrm{Hess}_{w_{\tau _1}}(J_{\tau _1}|_{\mathrm{\Phi }_{\tau _1}^1(z_{\tau _1})})\mathrm{ind}\mathrm{Hess}_{w_{\tau _0}}(J_{\tau _0}|_{\mathrm{\Phi }_{\tau _0}^1(z_{\tau _0})})`$ in terms of the family of Lagrangian subspaces $`\mathrm{\Lambda }(\tau )`$; that is to get a tool to follow the evolution of the Morse index under a continuous change of the parameters. This is indeed very useful since for some special values of the parameters the index could be known aโpriori. It concerns, in particular, optimal control problems with the parameter $`\tau =t_1t_0`$. If $`t_1t_0`$ is very small then sharpness of the Lagrangian point almost automatically implies the positivity or negativity of the Hessian.
First we discuss the finite-dimensional case: Theorem I.1 indicates that finite-dimensional approximations may already contain all essential information. Let $`Q_\tau `$ be a continuous family of quadratic forms defined on a finite-dimensional vector space. If $`\mathrm{ker}Q_\tau =0,\tau _0\tau \tau _1`$, then $`\mathrm{ind}Q_\tau `$ is constant on the segment $`[\tau _0,\tau _1]`$. This is why Lemma I.1 opens the way to follow evolution of the index in terms of the $``$-derivative: it locates values of the parameter where the index may change. Actually, $``$-derivative allows to evaluate this change as well; the increment of $`\mathrm{ind}Q_\tau `$ is computed via so called Maslov index of a family of Lagrangian subspaces. In order to define this index we have to recall some elementary facts about symplectic spaces.
Let $`\mathrm{\Sigma },\sigma `$ be a symplectic space, i.e. $`\mathrm{\Sigma }`$ is a $`2n`$-dimensional vector space and $`\sigma `$ be a nondegenerate anti-symmetric bilinear form on $`\mathrm{\Sigma }`$. The skew-orthogonal complement to the subspace $`\mathrm{\Gamma }\mathrm{\Sigma }`$ is the subspace $`\mathrm{\Gamma }^{\mathrm{}}=\{x\mathrm{\Sigma }:\sigma (x,\mathrm{\Gamma })=0\}`$. The nondegeneracy of $`\sigma `$ implies that $`dim\mathrm{\Gamma }^{\mathrm{}}=2ndim\mathrm{\Gamma }`$. A subspace $`\mathrm{\Gamma }`$ is isotropic if and only if $`\mathrm{\Gamma }^{\mathrm{}}\mathrm{\Gamma }`$; it is Lagrangian if and only if $`\mathrm{\Gamma }^{\mathrm{}}=\mathrm{\Gamma }`$.
Let $`\mathrm{\Pi }=span\{e_1,\mathrm{},e_n\}`$ be a lagrangian subspace of $`\mathrm{\Sigma }`$. Then there exist vectors $`f_1,\mathrm{},f_n\mathrm{\Sigma }`$ such that $`\sigma (e_i,f_j)=\delta _{ij}`$, where $`\delta _{ij}`$ is the Kronecker symbol. We show this using induction with respect to $`n`$. Skew-orthogonal complement to the space $`span\{e_1,\mathrm{},e_{n1}\}`$ contains an element $`f`$ which is not skew-orthogonal to $`e_n`$; we set $`f_n=\frac{1}{\sigma (e_n,f)}f`$. We have
$$span\{e_n,f_n\}span\{e_n,f_n\}^{\mathrm{}}=0$$
and the restriction of $`\sigma `$ to $`span\{e_n,f_n\}^{\mathrm{}}`$ is a nondegenerate bilinear form. Hence $`span\{e_n,f_n\}^{\mathrm{}}`$ is a $`2(n1)`$-dimensional symplectic space with a Lagrangian subspace $`span\{e_1,\mathrm{},e_{n1}\}`$. According to the induction assumption, there exist $`f_1,\mathrm{},f_{n1}`$ such that $`\sigma (e_i,f_j)=\delta _{ij}`$ and we are done.
Vectors $`e_1,\mathrm{},e_n,f_1,\mathrm{},f_n`$ form a basis of $`\mathrm{\Sigma }`$; in particular,$`\mathrm{\Delta }=span\{f_1,\mathrm{},f_n\}`$ is a transversal to $`\mathrm{\Pi }`$ Lagrangian subspace, $`\mathrm{\Sigma }=\mathrm{\Pi }\mathrm{\Delta }`$. If $`x_i=\underset{j=1}{\overset{n}{}}(\zeta _i^je_j+z_i^jf_j),i=1,2,`$ and $`\zeta _i=(\zeta _i^1,\mathrm{},\zeta _i^n)`$, $`z_i=(z_i^1,\mathrm{},z_i^n)^{}`$, then $`\sigma (x_1,x_2)=\zeta _1z_2\zeta _2z_1`$. The coordinates $`\zeta ,z`$ identify $`\mathrm{\Sigma }`$ with $`^n\times ^n`$; any transversal to $`\mathrm{\Delta }`$ $`n`$-dimensional subspace $`\mathrm{\Lambda }\mathrm{\Sigma }`$ has the following presentation in these coordinates:
$$\mathrm{\Lambda }=\{z^{},S_\mathrm{\Lambda }z):z^n\},$$
where $`S_\mathrm{\Lambda }`$ is an $`n\times n`$-matrix. The subspace $`\mathrm{\Lambda }`$ is Lagrangian if and only if $`S_\mathrm{\Lambda }^{}=S_\mathrm{\Lambda }`$. We have:
$$\mathrm{\Lambda }\mathrm{\Pi }=\{(z^{},0):z\mathrm{ker}S_\mathrm{\Lambda }\},$$
the subspace $`\mathrm{\Lambda }`$ is transversal to $`\mathrm{\Pi }`$ if and only if $`S_\mathrm{\Lambda }`$ is nondegenerate.
Thatโs time to introduce some notations. Let $`L(\mathrm{\Sigma })`$ be the set of all Lagrangian subspaces, a closed subset of the Grassmannian $`G_n(\mathrm{\Sigma })`$ of $`n`$-dimensional subspaces in $`\mathrm{\Sigma }`$. We set
$$\mathrm{\Delta }^{}=\{\mathrm{\Lambda }L(\mathrm{\Sigma }):\mathrm{\Lambda }\mathrm{\Delta }=0\},$$
an open subset of $`L(\mathrm{\Sigma })`$. The mapping $`\mathrm{\Lambda }S_\mathrm{\Lambda }`$ gives a regular parametrization of $`\mathrm{\Delta }^{}`$ by the $`n(n+1)/2`$-dimensional space of symmetric $`n\times n`$-matrices. Moreover, above calculations show that $`L(\mathrm{\Sigma })=\underset{\mathrm{\Delta }L(\mathrm{\Sigma })}{}\mathrm{\Delta }^{}`$. Hence $`L(\mathrm{\Sigma })`$ is a $`n(n+1)/2`$-dimensional submanifold of the Grassmannian $`G_n(\mathrm{\Sigma })`$ covered by coordinate charts $`\mathrm{\Delta }^{}`$. The manifold $`L(\mathrm{\Sigma })`$ is called Lagrange Grassmannian associated to the symplectic space $`\mathrm{\Sigma }`$. It is not hard to show that any coordinate chart $`\mathrm{\Delta }^{}`$ is everywhere dense in $`L(\mathrm{\Sigma })`$; our calculations give also a local parametrization of its complement.
Given $`\mathrm{\Pi }L(\mathrm{\Sigma })`$, the subset
$$_\mathrm{\Pi }=L(\mathrm{\Sigma })\mathrm{\Pi }^{}=\{\mathrm{\Lambda }L(\mathrm{\Sigma }):\mathrm{\Lambda }\mathrm{\Pi }0\}$$
is called the train of $`\mathrm{\Pi }`$. Let $`\mathrm{\Lambda }_0_\mathrm{\Pi },dim(\mathrm{\Lambda }_0\mathrm{\Pi })=k`$. Assume that $`\mathrm{\Delta }`$ is transversal to both $`\mathrm{\Lambda }_0`$ and $`\mathrm{\Pi }`$ (i.e. $`\mathrm{\Delta }\mathrm{\Lambda }_0^{}\mathrm{\Pi }^{}`$). The mapping $`\mathrm{\Lambda }S_\mathrm{\Lambda }`$ gives a regular parametrization of the neighborhood of $`\mathrm{\Lambda }_0`$ in $`_\mathrm{\Pi }`$ by a neighborhood of a corank $`k`$ matrix in the set of all degenerate symmetric $`n\times n`$-matrices. A basic perturbation theory for symmetric matrices now implies that a small enough neighborhood of $`\mathrm{\Lambda }_0`$ in $`_\mathrm{\Pi }`$ is diffeomorphic to the Cartesian product of a neighborhood of the origin of the cone of all degenerate symmetric $`k\times k`$-matrices and a $`(n(n+1)k(k+1))/2`$-dimensional smooth manifold (see \[1, Lemma 2.2\] for details). We see that $`_\mathrm{\Pi }`$ is not a smooth submanifold of $`L(\mathrm{\Sigma })`$ but a union of smooth strata, $`_\mathrm{\Pi }=\underset{k>0}{}_\mathrm{\Pi }^{(k)}`$, where $`_\mathrm{\Pi }^{(k)}=\{\mathrm{\Lambda }L(\mathrm{\Sigma }):dim(\mathrm{\Lambda }\mathrm{\Pi })=k\}`$ is a smooth submanifold of $`L(\mathrm{\Sigma })`$ of codimension $`k(k+1)/2`$.
Let $`\mathrm{\Lambda }(\tau ),\tau [t_0,t_1]`$ be a smooth family of Lagrangian subspaces (a smooth curve in $`L(\mathrm{\Sigma })`$) and $`\mathrm{\Lambda }(t_0),\mathrm{\Lambda }(t_1)\mathrm{\Pi }^{}`$. We are going to define the intersection number of $`\mathrm{\Lambda }()`$ and $`_\mathrm{\Pi }`$. It is called the Maslov index and is denoted $`\mu _\mathrm{\Pi }(\mathrm{\Lambda }())`$. Crucial property of this index is its homotopy invariance: given a homotopy $`\mathrm{\Lambda }^s()`$, $`s[t_0,t_1]`$ such that $`\mathrm{\Lambda }^s(t_0),\mathrm{\Lambda }^s(t_1)\mathrm{\Pi }^{}s[0,1]`$, we have $`\mu _\mathrm{\Pi }(\mathrm{\Lambda }^0())=\mu _\mathrm{\Pi }(\mathrm{\Lambda }^1())`$.
It is actually enough to define $`\mu _\mathrm{\Pi }(\mathrm{\Lambda }())`$ for the curves which have empty intersection with $`_\mathrm{\Pi }_\mathrm{\Pi }^{(1)}`$; the desired index would have a well-defined extension to other curves by continuity. Indeed, generic curves have empty intersection with $`_\mathrm{\Pi }_\mathrm{\Pi }^{(1)}`$ and, moreover, generic homotopy has empty intersection with $`_\mathrm{\Pi }_\mathrm{\Pi }^{(1)}`$ since any of submanifolds $`_\mathrm{\Pi }^{(k)},k=2,\mathrm{}n`$ has codimension greater or equal to 3 in $`L(\mathrm{\Sigma })`$. Putting any curve in general position by a small perturbation, we obtain the curve which bypasses $`_\mathrm{\Pi }_\mathrm{\Pi }^{(1)}`$, and the invariance with respect to generic homotopies of the Maslov index defined for generic curves would imply that the value of the index does not depend on the choice of a small perturbation.
What remains is to fix a โcoorientationโ of the smooth hypersurface $`_\mathrm{\Pi }^{(1)}`$ in $`L(\mathrm{\Sigma })`$, i. e. to indicate the โpositive and negative sidesโ of the hypersurface. As soon as we have a coorientation, we may compute $`\mu _\mathrm{\Pi }(\mathrm{\Lambda }())`$ for any curve $`\mathrm{\Lambda }()`$ which is transversal to $`_\mathrm{\Pi }^{(1)}`$ and has empty intersection with $`_\mathrm{\Pi }_\mathrm{\Pi }^{(1)}`$. Maslov index of $`\mathrm{\Lambda }()`$ is just the number of points where $`\mathrm{\Lambda }()`$ intersects $`_\mathrm{\Pi }^{(1)}`$ in the positive direction minus the number of points where this curve intersects $`_\mathrm{\Pi }^{(1)}`$ in the negative direction. Maslov index of any curve with endpoints out of $`_\mathrm{\Pi }`$ is defined by putting the curve in general position. Proof of the homotopy invariance is the same as for usual intersection number of a curve with a closed cooriented hypersurface (see, for instance, the nice elementary book by J. Milnor โTopology from the differential viewpointโ, 1965).
The coorientation is a byproduct of the following important structure on the tangent spaces to $`L(\mathrm{\Sigma })`$. It happens that any tangent vector to $`L(\mathrm{\Sigma })`$ at the point $`\mathrm{\Lambda }L(\mathrm{\Sigma })`$ can be naturally identified with a quadratic form on $`\mathrm{\Lambda }`$. Her we use the fact that $`\mathrm{\Lambda }`$ is not just a point in the Grassmannian but an $`n`$-dimensional linear space. To associate a quadratic form on $`\mathrm{\Lambda }`$ to the velocity $`\dot{\mathrm{\Lambda }}(t)T_{\mathrm{\Lambda }(t)}L(\mathrm{\Sigma })`$ of a smooth curve $`\mathrm{\Lambda }()`$ we proceed as follows: given $`x\mathrm{\Lambda }(t)`$ we take a smooth curve $`\tau x(\tau )`$ in $`\mathrm{\Sigma }`$ in such a way that $`x(\tau )\mathrm{\Lambda }(\tau ),\tau `$ and $`x(\tau )=x`$. Then we define a quadratic form $`\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}(t)(x),x\mathrm{\Lambda }(t)`$, by the formula $`\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}(t)(x)=\sigma (x,\dot{x}(t))`$.
The point is that $`\sigma (x,\dot{x}(t))`$ does not depend on the freedom in the choice of the curve $`\tau x(\tau )`$, although $`\dot{x}(t)`$ depends on this choice. Let us check the required property in the coordinates. We have $`x=(z^{},S_{\mathrm{\Lambda }(t)}z)`$ for some $`z^n`$ and $`x(\tau )=(z(\tau )^{},S_{\mathrm{\Lambda }(\tau )}z(\tau ))`$. Then
$$\sigma (x,\dot{x}(t))=z^{}(\dot{S}_{\mathrm{\Lambda }(t)}z+S_{\mathrm{\Lambda }(t)}\dot{z})\dot{z}^{}S_{\mathrm{\Lambda }(t)}z=z^{}\dot{S}_{\mathrm{\Lambda }(t)}z;$$
vector $`\dot{z}`$ does not show up. We have obtained a coordinate presentation of $`\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}(t)`$:
$$\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}(t)(z^{},S_{\mathrm{\Lambda }(t)}z)=z^{}\dot{S}_{\mathrm{\Lambda }(t)}z,$$
which implies that $`\dot{\mathrm{\Lambda }}\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}},\dot{\mathrm{\Lambda }}T_\mathrm{\Lambda }L(\mathrm{\Sigma })`$ is an isomorphism of $`T_\mathrm{\Lambda }L(\mathrm{\Sigma })`$ on the linear space of quadratic forms on $`\mathrm{\Lambda }`$.
We are now ready to define the coorientation of $`_\mathrm{\Pi }^{(1)}`$. Assume that $`\mathrm{\Lambda }(t)_\mathrm{\Pi }^{(1)}`$, i. e. $`\mathrm{\Lambda }(t)\mathrm{\Pi }=x`$ for some nonzero vector $`x\mathrm{\Sigma }`$. In coordinates, $`x=(z^{},0)`$, where $`x=\mathrm{ker}S_{\mathrm{\Lambda }(t)}`$. It is easy to see that $`\dot{\mathrm{\Lambda }}(t)`$ is transversal to $`_\mathrm{\Pi }^{(1)}`$ (i. e. $`\dot{S}_{\mathrm{\Lambda }(t)}`$ is transversal to the cone of degenerate symmetric matrices) if and only if $`\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}(t)(x)0`$ (i. e. $`z^{}\dot{S}_{\mathrm{\Lambda }(t)}z0`$). Vector $`x`$ is defined up to a scalar multiplier and $`\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}(t)(\alpha x)=\alpha ^2\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}(t)(x)`$ so that the sign of $`\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}(t)(x)`$ does not depend on the selection of $`x`$.
Definition. We say that $`\mathrm{\Lambda }()`$ intersects $`_\mathrm{\Pi }^{(1)}`$ at the point $`\mathrm{\Lambda }(t)`$ in the positive (negative) direction if $`\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}(t)(x)>0`$ ($`<0`$).
This definition completes the construction of the Maslov index. A weak point of the construction is the necessity to put the curve in general position in order to compute the intersection number. This does not look as an efficient way to thinks since putting the curve in general position is nothing else but a deliberate spoiling of a maybe nice and symmetric original object that makes even more involved the nontrivial problem of the localization of its intersection with $`_\mathrm{\Pi }`$. Fortunately, just the fact that Maslov index is homotopy invariant leads to a very simple and effective way of its computation without putting things in general position and without looking for the intersection points with $`_\mathrm{\Pi }`$.
###### Lemma I.2
Assume that $`\mathrm{\Pi }\mathrm{\Delta }=\mathrm{\Lambda }(\tau )\mathrm{\Delta }=0,\tau [t_0,t_1]`$. Then $`\mu _\mathrm{\Pi }(\mathrm{\Lambda }())=\mathrm{ind}S_{\mathrm{\Lambda }(t_0)}\mathrm{ind}S_{\mathrm{\Lambda }(t_1)}`$, where $`\mathrm{ind}S`$ is the Morse index of the quadratic form $`z^{}Sz,z^n`$.
Proof. The matrices $`S_{\mathrm{\Lambda }(t_0)}`$ and $`S_{\mathrm{\Lambda }(t_0)}`$ are nondegenerate since $`\mathrm{\Lambda }(t_0)\mathrm{\Pi }=\mathrm{\Lambda }(t_1)\mathrm{\Pi }=0`$ (we define the Maslov index only for the curves whose endpoins are out of $`_\mathrm{\Pi }`$). The set of nondegenerate quadratic forms with a prescribed value of the Morse index is a connected open subset of the linear space of all quadratic forms in $`n`$ variables. Hence homotopy invariance of the Maslov index implies that $`\mu _\mathrm{\Pi }(\mathrm{\Lambda }())`$ depends only on $`\mathrm{ind}S_{\mathrm{\Lambda }(t_0)}`$ and $`\mathrm{ind}S_{\mathrm{\Lambda }(t_1)}`$. It remains to compute $`\mu _\mathrm{\Pi }`$ of sample curves in $`\mathrm{\Delta }^{}`$, say, for segments of the curve $`\mathrm{\Lambda }()`$ such that
$$S_{\mathrm{\Lambda }(\tau )}=\left(\begin{array}{cccc}\tau 1& 0& \mathrm{}& 0\\ 0& \tau 2& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& \mathrm{}& \tau n\end{array}\right).$$
$`\mathrm{}`$
In general, given curve is not contained in the fixed coordinate neighborhood $`\mathrm{\Delta }^{}`$ but any curve can be divided into segments $`\mathrm{\Lambda }()|_{[\tau _i,\tau _{i+1}]},i=0,\mathrm{},l`$, in such a way that $`\mathrm{\Lambda }(\tau )\mathrm{\Delta }_i^{}\tau [\tau _i,\tau _{i+1}]`$, where $`\mathrm{\Delta }_i\mathrm{\Pi }^{},i=0,\mathrm{},l`$; then $`\mu _\mathrm{\Pi }(\mathrm{\Lambda }())=\underset{i}{}\mu _\mathrm{\Pi }\left(\mathrm{\Lambda }()|_{[\tau _i,\tau _{i+1}]}\right).`$
Lemma I.2 implies the following useful formula which is valid for the important class of monotone increasing curves in the Lagrange Grassmannian, i.e. the curves $`\mathrm{\Lambda }()`$ such that $`\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}(t)`$ are nonnegative quadratic forms: $`\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}(t)0,t`$.
###### Corollary I.1
Assume that $`\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}(\tau )0,\tau [t_0,t_1]`$ and $`\{\tau [t_0,t_1]:\mathrm{\Lambda }(\tau )\mathrm{\Pi }0\}`$ is a finite subset of $`(t_0,t_1)`$. Then
$$\mu _\mathrm{\Pi }\left(\mathrm{\Lambda }()\right)=\underset{\tau (t_0,t_1)}{}dim(\mathrm{\Lambda }(\tau )\mathrm{\Pi }).$$
$`\mathrm{}`$
Corollary I.1 can be also applied to the case of monotone decreasing curves defined by the inequality $`\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}(t)0,t`$; the change of parameter $`tt_0+t_1t`$ makes the curve monotone increasing and and change sign of the Maslov index.
Let me now recall that our interest to these symplectic playthings was motivated by the conditional minimum problems. As it was mentioned at the beginning of the section, we are going to apply this stuff to the case $`\mathrm{\Sigma }=T_\mathrm{}_\tau (T^{}M)`$, $`\mathrm{}_\tau T_{z_\tau }^{}M`$, $`\mathrm{\Pi }=T_\mathrm{}_\tau (T_{z_\tau }^{}M)`$, $`\mathrm{\Lambda }(\tau )=_{(\mathrm{}_\tau ,w_\tau )}(\overline{\mathrm{\Phi }}_\tau )`$, where $`z_\tau =\mathrm{\Phi }_\tau (w_\tau )`$. In this case, not only $`\mathrm{\Lambda }`$ but also $`\mathrm{\Pi }`$ and even symplectic space $`\mathrm{\Sigma }`$ depend on $`\tau `$. We thus have to define Maslov index in such situation. This is easy. We consider the bundle
$$\{(\xi ,\tau ):\xi T_\mathrm{}_\tau (T^{}M),t_0\tau t_1\}$$
$`(19)`$
over the segment $`[t_0,t_1]`$ induced from $`T(T^{}M)`$ by the mapping $`\tau \mathrm{}_\tau `$. Bundle (19) endowed with the symplectic structure and its subbundle
$$\{(\xi ,\tau ):\xi T_\mathrm{}_\tau (T_{z_\tau }^{}M)\}$$
are trivial as any bundle over a segment. More precisely, let $`t[t_0,t_1]`$, $`\mathrm{\Sigma }_t=T_\mathrm{}_t(T^{}M)`$, $`\mathrm{\Pi }_t=T_\mathrm{}_t(T_{z_t}^{}M)`$; then there exists a continuous with respect to $`\tau `$ family of linear symplectic mappings $`\mathrm{\Xi }_\tau :T_\mathrm{}_\tau (T^{}M)\mathrm{\Sigma }_t`$ such that $`\mathrm{\Xi }_\tau (T_\mathrm{}_\tau (T_{z_\tau }^{}M))=\mathrm{\Pi }_t,t_0\tau t_1,\mathrm{\Xi }_t=\mathrm{Id}`$. To any continuous family of Lagrangian subspaces $`\mathrm{\Lambda }(\tau )T_\mathrm{}_\tau (T^{}M)`$, where $`\mathrm{\Lambda }(t_i)\mathrm{\Pi }_{t_i}=0,i=0,1`$, we associate a curve $`\mathrm{\Xi }_.\mathrm{\Lambda }():\tau \mathrm{\Xi }_\tau \mathrm{\Lambda }(\tau )`$ in the Lagrange Grassmannian $`L(\mathrm{\Sigma }_t)`$ and set $`\mu (\mathrm{\Lambda }())\stackrel{def}{=}\mu _{\mathrm{\Pi }_t}(\mathrm{\Xi }_.\mathrm{\Lambda }())`$. Homotopy invariance of the Maslov index implies that $`\mu _{\mathrm{\Pi }_t}(\mathrm{\Xi }_.\mathrm{\Lambda }())`$ does not depend on the choice of $`t`$ and $`\mathrm{\Xi }_\tau `$.
###### Theorem I.2
Assume that $`dim๐ฒ<\mathrm{}`$,
$$\overline{\mathrm{\Phi }}_\tau =(J_\tau ,\mathrm{\Phi }_\tau ):๐ฒ\times M,\tau [t_0,t_1]$$
is a continuous one-parametric family of smooth mappings and $`(\mathrm{}_\tau ,w_\tau )`$ is a continuous family of their Lagrangian points such that $`\mathrm{}_\tau 0`$, $`w_\tau `$ is a regular point of $`\mathrm{\Phi }_\tau `$, and $`\mathrm{ker}D_{w_\tau }\mathrm{\Phi }_\tau \mathrm{ker}(D_{w_\tau }^2J_\tau \mathrm{}_\tau D_{w_\tau }^2\mathrm{\Phi }_\tau )=0`$, $`t_0\tau t_1`$. Let $`z_\tau =\mathrm{\Phi }(w_\tau )`$, $`\mathrm{\Lambda }(\tau )=_{(\mathrm{}_\tau ,w_\tau )}(\overline{\mathrm{\Phi }}_\tau )`$. If $`\mathrm{Hess}_{w_{t_i}}(J_{t_i}|_{\mathrm{\Phi }_{t_i}^1(z_{t_i})}),i=1,2`$, are nondegenerate, then
$$\mathrm{ind}\mathrm{Hess}_{w_{t_0}}(J_{t_0}|_{\mathrm{\Phi }_{t_0}^1(z_{t_0})})\mathrm{ind}\mathrm{Hess}_{w_{t_1}}(J_{t_1}|_{\mathrm{\Phi }_{t_1}^1(z_{t_1})})=\mu (\mathrm{\Lambda }()).$$
Remark. If $`\mathrm{}_\tau =0`$, then $`w_\tau `$ is a critical point of $`J_\tau `$ (without restriction to the level set of $`\mathrm{\Phi }_\tau `$). Theorem I.2 can be extended to this situation (with the same proof) if we additionally assume that $`\mathrm{ker}\mathrm{Hess}_{w_\tau }J_\tau =0`$ for any $`\tau `$ such that $`\mathrm{}_\tau =0`$.
Proof. We introduce simplified notations: $`A_\tau =D_{w_\tau }\mathrm{\Phi }_\tau `$, $`Q_\tau =D_{w_\tau }^2J_\tau \mathrm{}_\tau D_{w_\tau }^2\mathrm{\Phi }_\tau `$; the $``$-derivative $`_{(\mathrm{}_\tau ,w_\tau )}(\overline{\mathrm{\Phi }}_\tau )=\mathrm{\Lambda }(\tau )`$ is uniquely determined by the linear map $`A_\tau `$ and the symmetric bilinear form $`Q_\tau `$. Fix local coordinates in the neighborhoods of $`w_\tau `$ and $`z_\tau `$ and set:
$$\mathrm{\Lambda }(A,Q)=\{(\zeta ,Av):\zeta A+Q(v,)=0\}L(^n\times ^n);$$
then $`\mathrm{\Lambda }_\tau =\mathrm{\Lambda }(A_\tau ,Q_\tau )`$.
The assumption $`\mathrm{ker}A_\tau \mathrm{ker}Q_\tau =0`$ implies the smoothness of the mapping $`(A,Q)\mathrm{\Lambda }(A,Q)`$ for $`(A,Q)`$ close enough to $`(A_\tau ,Q_\tau )`$. Indeed, as it is shown in the proof of Proposition I.4, this assumption implies that the mapping $`left_\tau :(\zeta ,v)\zeta A_\tau +Q_\tau (v,)`$ is surjective. Hence the kernel of the mapping
$$(\zeta ,v)\zeta A+Q(v,)$$
$`(20)`$
smoothly depends on $`(A,Q)`$ for $`(A,Q)`$ close to $`(A_\tau ,Q_\tau )`$. On the other hand, $`\mathrm{\Lambda }(A,Q)`$ is the image of the mapping $`(\zeta ,v)(\zeta ,Av)`$ restricted to the kernel of map (20).
Now we have to disclose a secret which the attentive reader already knows and is perhaps indignant with our lightness: $`Q_\tau `$ is not a well-defined bilinear form on $`T_{w_\tau }๐ฒ`$, it essentially depends on the choice of local coordinates in $`M`$. What are well-defined is the mapping $`Q_\tau |_{\mathrm{ker}A_\tau }:\mathrm{ker}A_\tau T_{w_\tau }^{}๐ฒ`$ (check this by yourself or see \[3, Subsec. 2.3\]), the map $`A_\tau :T_{w_\tau }๐ฒT_{z_\tau }M`$ and, of course, the Lagrangian subspace $`\mathrm{\Lambda }(\tau )=_{(\mathrm{}_\tau ,w_\tau )}(\overline{\mathrm{\Phi }}_\tau )`$. By the way, the fact that $`Q_\tau |_{\mathrm{ker}A_\tau }`$ is well-defined guarantees that assumptions of Theorem I.2 do not depend on the coordinates choice.
Recall that any local coordinates $`\{z\}`$ on $`M`$ induce coordinates $`\{(\zeta ,z):\zeta ^n,z^n\}`$ on $`T^{}M`$ and $`T_z^{}M=\{(\zeta ,0):\zeta ^n\}`$ in the induced coordinates.
###### Lemma I.3
Given $`\widehat{z}M`$, $`\mathrm{}T_{\widehat{z}}^{}M\{0\}`$, and a Lagrangian subspace $`\mathrm{\Delta }T_{\mathrm{}}(T_{\widehat{z}}^{}M)^{}L(T_{\mathrm{}}(T^{}M))`$, there exist centered at $`\widehat{z}`$ local coordinates on $`M`$ such that $`\mathrm{\Delta }=\{(0,z):z^n\}`$ in the induced coordinates on $`T_{\mathrm{}}(T^{}M)`$.
Proof. Working in arbitrary local coordinates we have $`\mathrm{}=(\zeta _0,0)`$, $`\mathrm{\Delta }=\{(Sz,z):z^n\}`$, where $`S`$ is a symmetric matrix. In other words, $`\mathrm{\Delta }`$ is the tangent space at $`(\zeta _0,0)`$ to the graph of the differential of the function $`a(z)=\zeta _0z+\frac{1}{2}z^{}Sz`$. any smooth function with a nonzero differential can be locally made linear by a smooth change of variables. To prove the lemma it is enough to make a coordinates change which kills second derivative of the function $`a`$, for instance: $`zz+\frac{1}{2|\zeta _0|^2}(z^{}Sz)\zeta _0^{}.\mathrm{}`$
We continue the proof of Theorem I.2. Lemma I.3 gives us the way to take advantage of the fact that $`Q_\tau `$ depends on the choice of local coordinates in $`M`$. Indeed, bilinear form $`Q_\tau `$ is degenerate if and only if $`\mathrm{\Lambda }_\tau \{(0,z):z^n\}0`$; this immediately follows from the relation
$$\mathrm{\Lambda }_\tau =\{(\zeta ,A_\tau v):\zeta A_\tau +Q_\tau (v,)=0\}.$$
Given $`t[t_0,t_1]`$ take a transversal to $`T_\mathrm{}_t(T_{z_t}^{}M)`$ and $`\mathrm{\Lambda }(t)`$ Lagrangian subspace $`\mathrm{\Delta }_tT_\mathrm{}_t(T^{}M)`$ and centered at $`z_t`$ local coordinates in $`M`$ such that $`\mathrm{\Delta }_t=\{(0,z):z^n\}`$ in these coordinates. Then $`\mathrm{\Lambda }(\tau )`$ is transversal to $`\{(0,z):z^n\}`$ for all $`\tau `$ from a neighborhood $`O_t`$ of $`t`$ in $`[t_0,t_1]`$. Selecting an appropriate finite subcovering from the covering $`O_t,t[t_0,t_1]`$ of $`[t_0,t_1]`$ we can construct a subdivision $`t_0=\tau _0<\tau _1<\mathrm{}<\tau _k<\tau _{k+1}=t_1`$ of $`[t_0,t_1]`$ with the following property: $`i\{0,1,\mathrm{},k\}`$ the segment $`\{z_\tau :\tau [\tau _i,\tau _{i+1}]\}`$ of the curve $`z_\tau `$ is contained in a coordinate neighborhood $`๐ช^i`$ of $`M`$ such that $`\mathrm{\Lambda }_\tau \{(0,z):z^n\}=0\tau [\tau _i,\tau _{i+1}]`$ in the correspondent local coordinates.
We identify the form $`Q_\tau `$ with its symmetric matrix, i.e. $`Q_\tau (v_1,v_2)=v_1^{}Q_\tau v_2`$. Then $`Q_\tau `$ is a nondegenerate symmetric matrix and
$$\mathrm{\Lambda }(\tau )=\{(\zeta ,A_\tau Q_\tau ^1A_\tau ^{}\zeta ^{}\},\tau _i\tau \tau _{i+1}.$$
$`(21)`$
Now focus on the subspace $`\mathrm{\Lambda }(\tau _i)`$; it has a nontrivial intersection with $`\{(\zeta ,0):\zeta ^n\}=T_{\mathrm{}_{\tau _i}}(T_{z_{\tau _i}}^{}M)`$ if and only if the matrix $`A_{\tau _i}Q_{\tau _i}^1A_{\tau _i}^{}`$ is degenerate. This is the matrix of the restriction of the nondegenerate quadratic form $`vv^{}Q_{\tau _i}^1v`$ to the image of the linear map $`A_{\tau _i}^{}`$. Hence $`A_{\tau _i}Q_{\tau _i}^1A_{\tau _i}^{}`$ can be made nondegenerate by the arbitrary small perturbation of the map $`A_{\tau _i}:T_{w_{\tau _i}}๐ฒT_{z_{\tau _i}}M`$. Such perturbations can be realized simultaneously for $`i=1,\mathrm{},k`$<sup>2</sup><sup>2</sup>2We do not need to perturb $`A_{t_0}`$ and $`A_{t_{k+1}}`$: assumption of the theorem and Lemma I.1 guarantee the required nondegeneracy property. by passing to a continuous family $`\tau A_\tau ^{},t_0\tau t_1`$, arbitrary close and homotopic to the family $`\tau A_\tau `$. In fact, $`A_\tau ^{}`$ can be chosen equal to $`A_\tau `$ out of an arbitrarily small neighborhood of $`\{\tau _1,\mathrm{},\tau _k\}`$. Putting now $`A_\tau ^{}`$ instead of $`A_\tau `$ in the expression for $`\mathrm{\Lambda }(\tau )`$ we obtain a family of Lagrangian subspaces $`\mathrm{\Lambda }^{}(\tau )`$. This family is continuous (see the paragraph containing formula (20)) and homotopic to $`\mathrm{\Lambda }()`$. In particular, it has the same Maslov index as $`\mathrm{\Lambda }()`$. In other words, we can assume without lack of generality that $`\mathrm{\Lambda }(\tau _i)T_{\mathrm{}_{\tau _i}}(T_{z_{\tau _i}}^{}M)=0,i=0,1,\mathrm{},k+1`$. Then $`\mu (\mathrm{\Lambda }())=\underset{i=0}{\overset{k}{}}\mu \left(\mathrm{\Lambda }()|_{[\tau _i,\tau _{i+1}]}\right).`$ Moreover, it follows from (21) and Lemma I.2 that
$$\mu \left(\mathrm{\Lambda }()|_{[\tau _i,\tau _{i+1}]}\right)=\mathrm{ind}(A_{\tau _{i+1}}Q_{\tau _{i+1}}^1A_{\tau _{i+1}}^{})\mathrm{ind}(A_{\tau _i}Q_{\tau _i}^1A_{\tau _i}^{}).$$
Besides that, $`\mathrm{ind}Q_{\tau _i}=\mathrm{ind}Q_{\tau _{i+1}}`$ since $`Q_\tau `$ is nondegenerate for all $`\tau [\tau _i,\tau _{i+1}]`$ and continuously depends on $`\tau `$.
Recall that $`\mathrm{Hess}_{w_\tau }\left(J_\tau |_{\mathrm{\Phi }^1(z_\tau )}\right)=Q_\tau |_{\mathrm{ker}A_\tau }.`$ In order to complete proof of the theorem it remains to show that
$$\mathrm{ind}Q_\tau =\mathrm{ind}\left(Q_\tau |_{\mathrm{ker}A_\tau }\right)+\mathrm{ind}(A_\tau Q_\tau ^1A_\tau ^{})$$
$`(22)`$
for $`\tau =\tau _i,\tau _{i+1}`$.
Let us rearrange the second term in the right-hand side of (22). The change of variables $`v=Q_\tau ^1A_\tau ^{}z,z^n`$, implies: $`\mathrm{ind}\left(A_\tau Q_\tau ^1A_\tau ^{}\right)=\mathrm{ind}\left(Q_\tau |_{\{Q_\tau ^1A_\tau ^{}z:z^n\}}\right).`$ We have: $`Q_\tau (v,\mathrm{ker}A_\tau )=0`$ if and only if $`Q_\tau (v,)=z^{}A_\tau `$ for some $`z^n`$, i.e. $`v^{}Q_\tau =z^{}A_\tau `$, $`v=Q_\tau ^1A_\tau ^{}z`$. Hence the right-hand side of (22) takes the form
$$\mathrm{ind}Q_\tau =\mathrm{ind}\left(Q_\tau |_{\mathrm{ker}A_\tau }\right)+\mathrm{ind}\left(Q_\tau |_{\{v:Q_\tau (v,\mathrm{ker}A_\tau )=0\}}\right)$$
and $`Q_\tau |_{\{v:Q_\tau (v,\mathrm{ker}A_\tau )=0\}}`$ is a nondegenerate form for $`\tau =\tau _i,\tau _{i+1}`$. Now equality (22) is reduced to the following elementary fact of linear algebra: If $`Q`$ is a nondegenerate quadratic form on $`^m`$ and $`E^m`$ is a linear subspace, then $`\mathrm{ind}Q=\mathrm{ind}\left(Q|_E\right)+\mathrm{ind}\left(Q|_{E_Q^{}}\right)+dim(EE_Q^{}),`$ where $`E_Q^{}=\{v^m:Q(v,E)=0\}`$ and $`EE_Q^{}=\mathrm{ker}\left(Q|_E\right)=\mathrm{ker}\left(Q|_{E_Q^{}}\right).\mathrm{}`$
Remark. Maslov index $`\mu _\mathrm{\Pi }`$ is somehow more than just the intersection number with $`_\mathrm{\Pi }`$. It can be extended, in a rather natural way, to all continuous curves in the Lagrange Grassmannian including those whose endpoint belong to $`_\mathrm{\Pi }`$. This extension allows to get rid of the annoying nondegeneracy assumption for $`\mathrm{Hess}_{w_{t_i}}(J_{t_i}|_{\mathrm{\Phi }_{t_i}^1(z_{t_i})})`$ in the statement of Theorem I.2. In general, Maslov index computes 1/2 of the difference of the signatures of the Hessians which is equal to the difference of the Morse indices in the degenerate case (see for this approach).
### 7 Regular extremals
A combination of the finite-dimensional Theorem I.2 with the limiting procedure of Theorem I.1 and with homotopy invariance of the Maslov index allows to efficiently compute Morse indices of the Hessians for numerous infinite-dimensional problems. Here we restrict ourselves to the simplest case of a regular extremal of the optimal control problem.
We use notations and definitions of Sections 3, 4. Let $`h(\lambda ,u)`$ be the Hamiltonian of a smooth optimal control system and $`\lambda _t,t_0tt_1`$, be an extremal contained in the regular domain $`๐`$ of $`h`$. Then $`\lambda _t`$ is a solution of the Hamiltonian system $`\dot{\lambda }=\stackrel{}{H}(\lambda )`$, where $`H(\lambda )=h(\lambda ,\overline{u}(\lambda )),\frac{h}{u}h(\lambda ,\overline{u}(\lambda ))=0`$.
Let $`q(t)=\pi (\lambda _t),t_0,tt_1`$ be the extremal path. Recall that the pair $`(\lambda _{t_0},\lambda _t)`$ is a Lagrangian multiplier for the conditional minimum problem defined on an open subset of the space
$$M\times L_{\mathrm{}}([t_0,t_1],U)=\{(q_t),u()):qM,u()L_{\mathrm{}}([t_0,t_1],U)\},$$
where $`u()`$ is control and $`q_t`$ is the value at $`t`$ of the solution to the differential equation $`\dot{q}=f(q,u(\tau )),\tau [t_0,t_1]`$. In particular, $`F_t(q_t,u())=q_t`$. The cost is $`J_{t_0}^{t_1}(q_t,u())`$ and constraints are $`F_{t_0}(q_t,u())=q(0),q_t=q(t)`$.
Let us set $`J_t(u)=J_{t_0}^t(q(t),u()),\mathrm{\Phi }_t(u)=F_{t_0}(q(t),u())`$. A covector $`\lambda T^{}M`$ is a Lagrange multiplier for the problem $`(J_t,\mathrm{\Phi }_t)`$ if and only if there exists an extremal $`\widehat{\lambda }_\tau ,t_0\tau t`$, such that $`\lambda _{t_0}=\lambda ,\widehat{\lambda }_tT_{q(t)}^{}M`$. In particular, $`\lambda _{t_0}`$ is a Lagrange multiplier for the problem $`(J_t,\mathrm{\Phi }_t)`$ associated to the control $`u()=\overline{u}(\lambda _.)`$. Moreover, all sufficiently close to $`\lambda _{t_0}`$ Lagrange multipliers for this problem are values at $`t_0`$ of the solutions $`\lambda (\tau ),t_0\tau t`$ to the Hamiltonian system $`\dot{\lambda }=\stackrel{}{H}(\lambda )`$ with the boundary condition $`\lambda (t)T_{q(t)}^{}M`$.
Weโll use exponential notations for one-parametric groups of diffeomorphisms generated ordinary differential equations. In particular, $`e^{\tau \stackrel{}{H}}:T^{}MT^{}M,\tau `$, is a flow generated by the equation $`\dot{\lambda }=\stackrel{}{H}(\lambda )`$, so that $`\lambda (\tau ^{})=e^{(\tau ^{}\tau )\stackrel{}{H}}(\lambda (\tau ),\tau ,\tau ^{},`$ and Lagrange multipliers for the problem $`(J_t,\mathrm{\Phi }_t)`$ fill the $`n`$-dimensional submanifold $`e^{(t_0t)\stackrel{}{H}}\left(T_{q(t)}^{}M\right)`$.
We set $`\overline{\mathrm{\Phi }}_t=(J_t,\mathrm{\Phi }_t)`$; it is easy to see that the $``$-derivative $`_{(\lambda _{t_0},u)}(\overline{\mathrm{\Phi }}_t)`$ is the tangent space to $`e^{(t_0t)\stackrel{}{H}}\left(T_{q(t)}^{}M\right)`$, i.e. $`_{(\lambda _{t_0},u)}(\overline{\mathrm{\Phi }}_t)=e_{}^{(t_0t)\stackrel{}{H}}T_{\lambda _t}\left(T_{q(t)}^{}M\right)`$. Indeed, let us recall the construction of the $``$-derivative. First we linearize the equation for Lagrange multipliers at $`\lambda _{t_0}`$. Solutions of the linearized equation form an isotropic subspace $`_{(\lambda _{t_0},u)}^0(\overline{\mathrm{\Phi }}_t)`$ of the symplectic space $`T_{\lambda _{t_0}}(T^{}M)`$. If $`_{(\lambda _{t_0},u)}^0(\overline{\mathrm{\Phi }}_t)`$ is a Lagrangian subspace (i.e. $`dim_{(\lambda _{t_0},u)}^0(\overline{\mathrm{\Phi }}_t)=dimM`$), then $`_{(\lambda _{t_0}^0,u)}(\overline{\mathrm{\Phi }}_t)=_{(\lambda _{t_0},u)}(\overline{\mathrm{\Phi }}_t)`$, otherwise we need a limiting procedure to complete the Lagrangian subspace. In the case under consideration, $`_{(\lambda _{t_0},u)}^0(\overline{\mathrm{\Phi }}_t)=e_{}^{(t_0t)\stackrel{}{H}}T_{\lambda _t}\left(T_{q(t)}^{}M\right)`$ has a proper dimension and thus coincides with $`_{(\lambda _{t_0},u)}(\overline{\mathrm{\Phi }}_t)`$. We can check independently that $`e_{}^{(t_0t)\stackrel{}{H}}T_{\lambda _t}\left(T_{q(t)}^{}M\right)`$ is Lagrangian: indeed, $`T_{\lambda _t}\left(T_{q(t)}^{}M\right)`$ is Lagrangian and $`e_{}^{(t_0t)\stackrel{}{H}}:T_{\lambda _t}(T^{}M)T_{\lambda _{t_0}}(T^{}M)`$ is an isomorphism of symplectic spaces since Hamiltonian flows preserve the symplectic form.
So $`t_{(\lambda _{t_0},u)}(\overline{\mathrm{\Phi }}_t)`$ is a smooth curve in the Lagrange Grassmannian $`L\left(T_{\lambda _{t_0}}(T^{}M)\right)`$ and we can try to compute Morse index of
$$\mathrm{Hess}_u\left(J_{t_1}|_{\mathrm{\Phi }_{t_1}^1(q(t_0))}\right)=\mathrm{Hess}_u\left(J_{t_0}^{t_1}|_{F_{t_0}^1(q(t_0))F_{t_1}^1(q(t_1))}\right)$$
via the Maslov index of this curve. Of course, such a computation has no sense if the index is infinite.
###### Proposition I.5
(Legendre condition) If quadratic form $`\frac{^2h}{u^2}(\lambda _t,u(t))`$ is negative definite for any $`t[t_0,t_1]`$, then $`\mathrm{ind}\mathrm{Hess}_u\left(J_{t_1}|_{\mathrm{\Phi }_{t_1}^1(q(t_0))}\right)<\mathrm{}`$ and $`\mathrm{Hess}_u\left(J_t|_{\mathrm{\Phi }_t^1(q(t_0))}\right)`$ is positive definite for any $`t`$ sufficiently close to (and strictly greater than) $`t_0`$. If $`\frac{^2h}{u^2}(\lambda _t,u(t))0`$ for some $`t[t_0,t_1]`$, then $`\mathrm{ind}\mathrm{Hess}_u\left(J_{t_1}|_{\mathrm{\Phi }_{t_1}^1(q(t_0))}\right)=\mathrm{}`$.
We do not give here the proof of this well-known result; you can find it in many sources (see, for instance, the textbook ). It is based on the fact that $`\frac{^2h}{u^2}(\lambda _t,u(t))=\lambda (\frac{^2f}{u^2}(q(t),u(t)))\frac{^2\phi }{u^2}(q(t),u(t))`$ is the infinitesimal (for the โinfinitesimally small intervalโ at $`t`$) version of $`\lambda _{t_0}D_u^2\mathrm{\Phi }_{t_1}D_u^2J_{t_1}`$ while $`\mathrm{Hess}_u(J_{t_1}|_{\mathrm{\Phi }_{t_1}^1(q(t_0))})=(D_u^2J_{t_1}\lambda _{t_0}D_w^2\mathrm{\Phi }_{t_1})|_{\mathrm{ker}D_u\mathrm{\Phi }_{t_1}}`$.
Next theorem shows that in the โregularโ infinite dimensional situation of this section we may compute the Morse index similarly to the finite dimensional case. The proof of the theorem requires some information about second variation of optimal control problems which is out of the scope of these notes. The required information can be found in Chapters 20, 21 of . Basically, it implies that finite dimensional arguments used in the proof of Theorem I.2 are legal also in our infinite dimensional case.
We set: $`\mathrm{\Lambda }(t)=e_{}^{(t_0t)\stackrel{}{H}}T_{\lambda _t}\left(T_{q(t)}^{}M\right)`$.
###### Theorem I.3
Assume that $`\frac{^2h}{u^2}(\lambda _t,u(t))`$ is a negative definite quadratic form and $`u`$ is a regular point of $`\mathrm{\Phi }_t,t(t_0,t_1].`$ Then:
* The form $`\mathrm{Hess}_u\left(J_{t_1}|_{\mathrm{\Phi }_{t_1}^1(q(t_0))}\right)`$ is degenerate if and only if
$`\mathrm{\Lambda }(t_1)\mathrm{\Lambda }(t_0)0`$;
* If $`\mathrm{\Lambda }(t_1)\mathrm{\Lambda }(t_0)=0`$, then there exists $`\overline{t}>t_0`$ such that
$$\mathrm{ind}\mathrm{Hess}_u\left(J_{t_1}|_{\mathrm{\Phi }_{t_1}^1(q(t_0))}\right)=\mu \left(\mathrm{\Lambda }()|_{[\tau ,t_1]}\right),\tau (t_0,\overline{t}).$$
$`\mathrm{}`$
Note that Legendre condition implies monotonicity of the curve $`\mathrm{\Lambda }()`$; this property simplifies the evaluation of the Maslov index. Fix some local coordinates in $`M`$ so that $`T^{}M\{(p,q)^n\times ^n\}`$.
###### Lemma I.4
Quadratic form $`\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}(t)`$ is equivalent (with respect to a linear change of variables) to the form $`\frac{^2H}{p^2}(\lambda _t)=\frac{\overline{u}}{p}^{}\frac{^2h}{u^2}(\lambda _t,\overline{u}(\lambda _t))\frac{\overline{u}}{p}`$.
Proof. Equality $`\frac{^2H}{p^2}=\frac{\overline{u}}{p}^{}\frac{^2h}{u^2}\frac{\overline{u}}{p}`$ is an easy corollary of the identities $`H(p,q)=h(p,q,\overline{u}(p,q)),\frac{h}{u}|_{u=\overline{u}(p,q)}=0`$. Indeed, $`\frac{^2H}{p^2}=2\frac{^2h}{up}\frac{\overline{u}}{p}+\frac{\overline{u}}{p}^{}\frac{^2h}{u^2}\frac{\overline{u}}{p}`$ and $`\frac{}{p}\left(\frac{h}{u}\right)=\frac{^2h}{pu}+\frac{^2h}{u^2}\frac{\overline{u}}{p}=0`$. Further, we have:
$$\frac{d}{dt}\mathrm{\Lambda }(t)=\frac{d}{dt}e_{}^{(t_0t)\stackrel{}{H}}T_{\lambda _t}\left(T_{q(t)}^{}M\right)=e_{}^{(t_0t)\stackrel{}{H}}\frac{d}{d\epsilon }|_{\epsilon =0}e_{}^{\epsilon \stackrel{}{H}}T_{\lambda _{t+\epsilon }}\left(T_{q(t+\epsilon )}^{}M\right).$$
Set $`\mathrm{\Delta }(\epsilon )=e_{}^{\epsilon \stackrel{}{H}}T_{\lambda _{t+\epsilon }}\left(T_{q(t+\epsilon )}^{}M\right)L\left(T_{\lambda (t)}(T^{}M)\right)`$. It is enough to prove that $`\underset{ยฏ}{\dot{\mathrm{\Delta }}(0)}`$ is equivalent to $`\frac{^2H}{p^2}(\lambda _t)`$. Indeed, $`\dot{\mathrm{\Lambda }}(t)=e_{}^{(t_0t)\stackrel{}{H}}T_{\lambda _t}\dot{\mathrm{\Delta }}(0)`$, where
$$e_{}^{(t_0t)\stackrel{}{H}}:T_{\lambda _t}(T^{}M)T_{\lambda _{t_0}}(T^{}M)$$
is a symplectic isomorphism. The association of the quadratic form $`\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}(t)`$ on the subspace $`\mathrm{\Lambda }(t)`$ to the tangent vector $`\dot{\mathrm{\Lambda }}(t)L\left(T_{\lambda _{t_0}}(T^{}M)\right)`$ is intrinsic, i.e. depends only on the symplectic structure on $`(T_{\lambda _{t_0}}(T^{}M)`$. Hence $`\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Delta }}}(0)(\xi )=\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}(t)\left(e_{}^{(t_0t)\stackrel{}{H}}\xi \right)`$, $`\xi \mathrm{\Delta }(0)=T_{\lambda _t}\left(T_{q(t)}^{}M\right)`$.
What remains, is to compute $`\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Delta }}}(0)`$; we do it in coordinates. We have:
$$\mathrm{\Delta }(\epsilon )=\{(\xi (\epsilon ),\eta (\epsilon )):\begin{array}{ccc}\hfill \dot{\xi }(\tau )& =& \xi \frac{^2H}{pq}(\lambda _{t\tau })+\eta ^{}\frac{^2H}{q^2}(\lambda _{t\tau }),\hfill \\ \hfill \dot{\eta }(\tau )& =& \frac{^2H}{p^2}(\lambda _{t\tau })\xi ^{}\frac{^2H}{qp}(\lambda _{t\tau })\eta ,\hfill \end{array}\genfrac{}{}{0pt}{}{\xi (0)^n}{\eta (0)=0}\},$$
$$\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Delta }}}(0)(\xi (0))=\sigma ((\xi (0),0),(\dot{\xi }(0),\dot{\eta }(0)))=\xi (0)\dot{\eta }(0)=\xi (0)\frac{^2H}{p^2}(\lambda _t)\xi (0)^{}.$$
$`\mathrm{}`$
Now combining Lemma I.4 with Theorem I.3 and Corollary I.1 we obtain the following version of the classical โMorse formulaโ
###### Corollary I.2
Under conditions of Theorem I.3, if $`\{\tau (t_0,t_1]:\mathrm{\Lambda }(\tau )\mathrm{\Lambda }(t_0)0\}`$ is a finite subset of $`(t_0,t_1)`$, then
$$\mathrm{ind}\mathrm{Hess}J_{t_1}|_{\mathrm{\Phi }_{t_1}^1(q(t_0))}=\underset{\tau (t_0,t_1)}{}dim(\mathrm{\Lambda }(\tau )\mathrm{\Lambda }(t_0)).$$
## Part II Geometry of Jacobi curves
### 8 Jacobi curves
Computation of the $``$-derivative for regular extremals in the last section has led us to the construction of curves in the Lagrange Grassmannians which works for all Hamiltonian systems on the cotangent bundles, independently on any optimal control problem. Set $`\mathrm{\Delta }_\lambda =T_\lambda (T_q^{}M)`$, where $`\lambda T_q^{}M,qM`$. The curve $`\tau e_{}^{\tau \stackrel{}{H}}\mathrm{\Delta }_{e^{\tau \stackrel{}{H}}(\lambda )}`$ in the Lagrange Grassmannian $`L\left(T_\lambda (T^{}M)\right)`$ is the result of the action of the flow $`e^{t\stackrel{}{H}}`$ on the vector distribution $`\{\mathrm{\Delta }_\lambda \}_{\lambda T^{}M}`$. Now we are going to study differential geometry of these curves; their geometry will provide us with a canonical connection on $`T^{}M`$ associated with the Hamiltonian system and with curvature-type invariants. All that gives a far going generalization (and a dynamical interpretation) of classical objects from Riemannian geometry.
In fact, construction of the basic invariants does not need symplectic structure and the Hamiltonian nature of the flow, we may deal with more or less arbitrary pairs (vector field, rank $`n`$ distribution) on a $`2n`$-dimensional manifold $`N`$. The resulting curves belong to the usual Grassmannian of all $`n`$-dimensional subspaces in the $`2n`$-dimensional one. We plan to work for some time in this more general situation and then come back to the symplectic framework.
In these notes we mainly deal with the case of involutive distributions (i.e. with $`n`$-foliations) just because our main motivation and applications satisfy this condition. The reader can easily recover more general definitions and construction by himself.
So we consider a $`2n`$-dimensional smooth manifold $`N`$ endowed with a smooth foliation of rank $`n`$. Let $`zN`$, by $`E_z`$ we denote the passing through $`z`$ leaf of the foliation; then $`E_z`$ is an $`n`$-dimensional submanifold of $`N`$. Point $`z`$ has a coordinate neighborhood $`O_z`$ such that the restriction of the foliation to $`O_z`$ is a (trivial) fiber bundle and the fibers $`E_z^{}^{loc},z^{}O_z,`$ of this fiber bundle are connected components of $`E_z^{}O_z`$. Moreover, there exists a diffeomorphism $`O_z^n\times ^n`$, where $`^n\times \{y\},y^n,`$ are identified with the fibers so that both the typical fiber and the base are diffeomorphic to $`^n`$. We denote by $`O_z/E^{loc}`$ the base of this fiber bundle and by $`\pi :O_zO_z/E^{loc}`$ the canonical projection.
Let $`\zeta `$ be a smooth vector field on $`N`$. Then $`z^{}\pi _{}\zeta (z^{})`$, $`z^{}E_z^{loc}`$ is a smooth mapping of $`E_z^{loc}`$ to $`T_{\pi (z)}(O_z/E^{loc})`$. We denote the last mapping by $`\mathrm{\Pi }_z(\zeta ):E_z^{loc}T_{\pi (z)}(O_z/E^{loc})`$.
Definition. We call $`\zeta `$ a lifting field if $`\mathrm{\Pi }_z(\zeta )`$ is a constant mapping $`zN`$; The field $`\zeta `$ is called regular if $`\mathrm{\Pi }_z(\zeta )`$ is a submersion, $`zN`$.
The flow generated by the lifting field maps leaves of the foliation in the leaves, in other words it is leaves-wise. On the contrary, the flow generated by the regular field โsmearsโ the fibers over $`O_z/E^{loc}`$; basic examples are second order differential equations on a manifold $`M`$ treated as the vector fields on the tangent bundle $`TM=N`$.
Let us write things in coordinates: We fix local coordinates acting in the domain $`ON`$, which turn the foliation into the Cartesian product of vector spaces: $`O\{(x,y):x,y^n\}`$, $`\pi :(x,y)y`$. Then vector field $`\zeta `$ takes the form $`\zeta =\underset{i=1}{\overset{n}{}}\left(a^i\frac{}{x_i}+b^i\frac{}{y_i}\right)`$, where $`a^i,b^i`$ are smooth functions on $`^n\times ^n`$. The coordinate representation of the map $`\mathrm{\Pi }_z`$ is: $`\mathrm{\Pi }_{(x,y)}:x(b^1(x,y),\mathrm{},b^n(x,y))^{}`$. Field $`\zeta `$ is regular if and only if $`\mathrm{\Pi }_{(x,y)}`$ are submersions; in other words, if and only if $`\left(\frac{b^i}{x_j}\right)_{i,j=1}^n`$ is a nondegenerate matrix. Field $`\zeta `$ is lifting if and only if $`\frac{b^i}{x_j}0,i,j=1,\mathrm{},n`$.
Now turn back to the coordinate free setting. The fibers $`E_z`$, $`zN`$ are integral manifolds of the involutive distribution $`=\{T_zE_z:zN\}`$. Given a vector field $`\zeta `$ on $`N`$, the (local) flow $`e^{t\zeta }`$ generated by $`\zeta `$, and $`zN`$ we define the family of subspaces
$$J_z(t)=\left(e^{t\zeta }\right)_{}|_zT_zN.$$
In other words, $`J_z(t)=\left(e^{t\zeta }\right)_{}T_{e^{t\zeta }(z)}E_{e^{t\zeta }(z)}`$, $`J_z(0)=T_zE_z`$.
$`J_x(t)`$ is an $`n`$-dimensional subspace of $`T_zN`$, i.e. an element of the Grassmannian $`G_n(T_zN)`$. We thus have (the germ of) a curve $`tJ_z(t)`$ in $`G_n(T_zN)`$ which is called a Jacobi curve.
Definition. We say that field $`\zeta `$ is k-ample for an interger $`k`$ if $`zN`$ and for any curve $`t\widehat{J}_z(t)`$ in $`G_n(T_zN)`$ with the same $`k`$-jet as $`J_z(t)`$ we have $`\widehat{J}_z(0)\widehat{J}_z(t)=0`$ for all $`t`$ close enough but not equal to 0. The field is called ample if it is $`k`$-ample for some $`k`$.
It is easy to show that a field is 1-ample if and only if it is regular.
### 9 The cross-ratio
Let $`\mathrm{\Sigma }`$ be a $`2n`$-dimensional vector space, $`v_0,v_1G_n(\mathrm{\Sigma }),v_0v_1=0`$. Than $`\mathrm{\Sigma }=v_0+v_1`$. We denote by $`\pi _{v_0v_1}:\mathrm{\Sigma }v_1`$ the projector of $`\mathrm{\Sigma }`$ onto $`v_1`$ parallel to $`v_0`$. In other words, $`\pi _{v_0v_1}`$ is a linear operator on $`\mathrm{\Sigma }`$ such that $`\pi _{v_0v_1}|_{v_0}=0`$, $`\pi _{v_0v_1}|_{v_1}=\text{id}`$. Surely, there is a one-to-one correspondence between pairs of transversal $`n`$-dimensional subspaces of $`\mathrm{\Sigma }`$ and rank $`n`$ projectors in $`\text{gl}(\mathrm{\Sigma })`$.
###### Lemma II.1
Let $`v_0G_n(\mathrm{\Sigma })`$; we set $`v_0^{}=\{vG_n(\mathrm{\Sigma }):vv_0=0\}`$, an open dense subset of $`G_n(\mathrm{\Sigma })`$. Then $`\{\pi _{vv_0}:vv_0^{}\}`$ is an affine subspace of $`\text{gl}(\mathrm{\Sigma })`$.
Indeed, any operator of the form $`\alpha \pi _{vv_0}+(1\alpha )\pi _{wv_0}`$, where $`\alpha `$, takes values in $`v_0`$ and its restriction to $`v_0`$ is the identity operator. Hence $`\alpha \pi _{vv_0}+(1\alpha )\pi _{wv_0}`$ is the projector of $`\mathrm{\Sigma }`$ onto $`v_0`$ along some subspace.
The mapping $`v\pi _{vv_0}`$ thus serves as a local coordinate chart on $`G_n(\mathrm{\Sigma })`$. These charts indexed by $`v_0`$ form a natural atlas on $`G_n(\mathrm{\Sigma })`$.
Projectors $`\pi _{vw}`$ satisfy the following basic relations:<sup>3</sup><sup>3</sup>3Numbering of formulas is separate in each of two parts of the paper
$$\pi _{v_0v_1}+\pi _{v_1v_0}=id,\pi _{v_0v_2}\pi _{v_1v_2}=\pi _{v_1v_2},\pi _{v_0v_1}\pi _{v_0v_2}=\pi _{v_0v_1},$$
$`(1)`$
where $`v_iG_n(\mathrm{\Sigma }),v_iv_j=0`$ for $`ij`$. If $`n=1`$, then $`G_n(\mathrm{\Sigma })`$ is just the projective line $`^1`$; basic geometry of $`G_n(\mathrm{\Sigma })`$ is somehow similar to geometry of the projective line for arbitrary $`n`$ as well. The group $`\text{GL}(\mathrm{\Sigma })`$ acts transitively on $`G_n(\mathrm{\Sigma })`$. Let us consider its standard action on $`(k+1)`$-tuples of points in $`G_n(\mathrm{\Sigma })`$:
$$A(v_0,\mathrm{},v_k)\stackrel{def}{=}(Av_0,\mathrm{},Av_k),A\text{GL}(\mathrm{\Sigma }),v_iG_n(\mathrm{\Sigma }).$$
It is an easy exercise to check that the only invariants of a triple $`(v_0,v_1,v_2)`$ of points of $`G_n(\mathrm{\Sigma })`$ for such an action are dimensions of the intersections: $`dim(v_iv_j),0i2`$, and $`dim(v_0v_1v_2)`$. Quadruples of points possess a more interesting invariant: a multidimensional version of the classical cross-ratio.
Definition. Let $`v_iG_n(\mathrm{\Sigma }),i=0,1,2,3`$, and $`v_0v_1=v_2v_3=0.`$ The cross-ratio of $`v_i`$ is the operator $`[v_0,v_1,v_2,v_3]\text{gl}(v_1)`$ defined by the formula:
$$[v_0,v_1,v_2,v_3]=\pi _{v_0v_1}\pi _{v_2v_3}|_{v_1}.$$
Remark. We do not lose information when restrict the product $`\pi _{v_0v_1}\pi _{v_2v_3}`$ to $`v_1`$; indeed, this product takes values in $`v_1`$ and its kernel contains $`v_0`$.
For $`n=1`$, $`v_1`$ is a line and $`[v_0,v_1,v_2,v_3]`$ is a real number. For general $`n`$, the Jordan form of the operator provides numerical invariants of the quadruple $`v_i,i=0,1,2,3`$.
We will mainly use an infinitesimal version of the cross-ratio that is an invariant $`[\xi _0,\xi _1]\text{gl}(v_1)`$ of a pair of tangent vectors $`\xi _iT_{v_i}G_n(\mathrm{\Sigma }),i=0,1,`$ where $`v_0v_1=0`$. Let $`\gamma _i(t)`$ be curves in $`G_n(\mathrm{\Sigma })`$ such that $`\gamma _i(0)=v_i,\frac{d}{dt}\gamma _i(t)|_{t=0}=\xi _i`$, $`i=0,1`$. Then the cross-ratio: $`[\gamma _0(t),\gamma _1(0),\gamma _0(\tau ),\gamma _1(\theta )]`$ is a well defined operator on $`v_1=\gamma _1(0)`$ for all $`t,\tau ,\theta `$ close enough to 0. Moreover, it follows from (1) that $`[\gamma _0(t),\gamma _1(0),\gamma _0(0),\gamma _1(0)]=`$$`[\gamma _0(0),\gamma _1(0),\gamma _0(t),\gamma _1(0)]=[\gamma _0(0),\gamma _1(0),\gamma _0(0),\gamma _1(t)]=id`$. We set
$$[\xi _0,\xi _1]=\frac{^2}{t\tau }[\gamma _0(t),\gamma _1(0),\gamma _0(0),\gamma _1(\tau )]|_{v_1}|_{t=\tau =0}$$
$`(2)`$
It is easy to check that the right-hand side of (2) depends only on $`\xi _0,\xi _1`$ and that $`(\xi _0,\xi _1)[\xi _0,\xi _1]`$ is a bilinear mapping from $`T_{v_0}G_n(\mathrm{\Sigma })\times T_{v_1}G_n(\mathrm{\Sigma })`$ onto $`gl(v_1)`$.
###### Lemma II.2
Let $`v_0,v_1G_n(\mathrm{\Sigma }),v_0v_1=0,\xi _iT_{v_i}G_n(\mathrm{\Sigma }),and\xi _i=\frac{d}{dt}\gamma _i(t)|_{t=0},i=0,1`$. Then $`[\xi _0,\xi _1]=\frac{^2}{t\tau }\pi _{\gamma _1(t)\gamma _0(\tau )}|_{v_1}|_{t=\tau =0}`$ and $`v_1,v_0`$ are invariant subspaces of the operator $`\frac{^2}{t\tau }\pi _{\gamma _1(t)\gamma _0(\tau )}|_{v_1}|_{t=\tau =0}`$.
Proof. According to the definition, $`[\xi _0,\xi _1]=\frac{^2}{t\tau }(\pi _{\gamma _0(t)\gamma _1(0)}\pi _{\gamma _0(0)\gamma _1(\tau )})|_{v_1}|_{t=\tau =0}.`$ The differentiation of the identities $`\pi _{\gamma _0(t)\gamma _1(0)}\pi _{\gamma _0(t)\gamma _1(\tau )}=\pi _{\gamma _0(t)\gamma _1(0)},`$$`\pi _{\gamma _0(t)\gamma _1(\tau )}\pi _{\gamma _0(0)\gamma _1(\tau )}=\pi _{\gamma _0(0)\gamma _1(\tau )}`$ gives the equalities:
$$\frac{^2}{t\tau }(\pi _{\gamma _0(t)\gamma _1(0)}\pi _{\gamma _0(0)\gamma _1(\tau )})|_{t=\tau =0}=\pi _{v_0v_1}\frac{^2}{t\tau }\pi _{\gamma _0(t)\gamma _1(\tau )}|_{t=\tau =0}$$
$$=\frac{^2}{t\tau }\pi _{\gamma _0(t)\gamma _1(\tau )}|_{t=\tau =0}\pi _{v_0v_1}.$$
It remains to mention that $`\frac{^2}{t\tau }\pi _{\gamma _1(t)\gamma _0(\tau )}=\frac{^2}{t\tau }\pi _{\gamma _0(\tau )\gamma _1(t)}`$. $`\mathrm{}`$
### 10 Coordinate setting
Given $`v_iG_n(\mathrm{\Sigma })`$, $`i=0,1,2,3`$, we coordinatize $`\mathrm{\Sigma }=^n\times ^n=\{(x,y):x^n,y^n\}`$ in such a way that $`v_i\{(0,y):y^n\}=0`$. Then there exist $`n\times n`$-matrices $`S_i`$ such that
$$v_i=\{(x,S_ix):x^n\},i=0,1,2,3.$$
$`(3)`$
The relation $`v_iv_j=0`$ is equivalent to $`det(S_iS_j)0`$. If $`S_0=0`$, then the projector $`\pi _{v_0v_1}`$ is represented by the $`2n\times 2n`$-matrix $`\left(\begin{array}{cc}0& S_1^1\\ 0& I\end{array}\right).`$ In general, we have
$$\pi _{v_0v_1}=\left(\begin{array}{cc}S_{01}^1S_0& S_{01}^1\\ S_1S_{01}^1S_0& S_1S_{01}^1\end{array}\right),$$
where $`S_{01}=S_0S_1`$. Relation (3) provides coordinates $`\{x\}`$ on the spaces $`v_i`$. In these coordinates, the operator $`[v_0,v_1,v_2,v_3]`$ on $`v_1`$ is represented by the matrix:
$$[v_0,v_1,v_2,v_3]=S_{10}^1S_{03}S_{32}^1S_{21},$$
where $`S_{ij}=S_iS_j`$.
We now compute the coordinate representation of the infinitesimal cross-ratio. Let $`\gamma _0(t)=\{(x,S_tx):x^n\}`$, $`\gamma _1(t)=\{(x,S_{1+t}x):x^n\}`$ so that $`\xi _i=\frac{d}{dt}\gamma _i(t)|_{t=0}`$ is represented by the matrix $`\dot{S}_i=\frac{d}{dt}S_t|_{t=i},i=0,1.`$ Then $`[\xi _0,\xi _1]`$ is represented by the matrix
$$\frac{^2}{t\tau }S_{1t}^1S_{t\tau }S_{\tau 0}^1S_{01}|_{\frac{t=0}{\tau =1}}=\frac{}{t}S_{1t}^1\dot{S}_1|_{t=0}=S_{01}^1\dot{S}_0S_{01}^1\dot{S}_1.$$
So
$$[\xi _0,\xi _1]=S_{01}^1\dot{S}_0S_{01}^1\dot{S}_1.$$
$`(4)`$
There is a canonical isomorphism $`T_{v_0}G_n(\mathrm{\Sigma })\text{Hom}(v_0,\mathrm{\Sigma }/v_0)`$; it is defined as follows. Let $`\xi T_{v_0}G_n(\mathrm{\Sigma }),\xi =\frac{d}{dt}\gamma (t)|_{t=0}`$, and $`z_0v_0`$. Take a smooth curve $`z(t)\gamma (t)`$ such that $`z(0)=z_0`$. Then the residue class $`(\dot{z}(0)+v_0)\mathrm{\Sigma }/v_0`$ depends on $`\xi `$ and $`z_0`$ rather than on a particular choice of $`\gamma (t)`$ and $`z(t)`$. Indeed, let $`\gamma ^{}(t)`$ be another curve in $`G_n(\mathrm{\Sigma })`$ whose velocity at $`t=0`$ equals $`\xi `$. Take some smooth with respect to $`t`$ bases of $`\gamma (t)`$ and $`\gamma ^{}(t)`$: $`\gamma (t)=span\{e_1(t),\mathrm{},e_n(t)\},\gamma ^{}(t)=span\{e_1^{}(t),\mathrm{},e_n^{}(t)\}`$, where $`e_i(0)=e_i^{}(0),i=1,\mathrm{},n`$; then $`\left(\dot{e}_i(0)\dot{e}_i^{}(0)\right)v_0,i=1,\mathrm{},n`$. Let $`z(t)=\underset{i=1}{\overset{n}{}}\alpha _i(t)e_i(t),z^{}(t)=\underset{i=1}{\overset{n}{}}\alpha _i^{}(t)e_i^{}(t)`$, where $`\alpha _i(0)=\alpha _i^{}(0)`$. We have:
$$\dot{z}(0)\dot{z}^{}(0)=\underset{i=1}{\overset{n}{}}\left((\dot{\alpha }_i(0)\dot{\alpha }_i^{}(0))e_i(0)+\alpha _i^{}(0)(\dot{e}_i(0)\dot{e}_i^{}(0))\right)v_0,$$
i.e. $`\dot{z}(0)+v_0=\dot{z}^{}(0)+v_0`$.
We associate to $`\xi `$ the mapping $`\overline{\xi }:v_0\mathrm{\Sigma }/v_0`$ defined by the formula $`\overline{\xi }z_0=\dot{z}(0)+v_0`$. The fact that $`\xi \overline{\xi }`$ is an isomorphism of the linear spaces $`T_{v_0}G_n(\mathrm{\Sigma })`$ and $`\text{Hom}(v_0,\mathrm{\Sigma }/v_0)`$ can be easily checked in coordinates. The matrices $`\dot{S}_i`$ above are actually coordinate presentations of $`\overline{\xi }_i,i=0,1`$.
The standard action of the group $`\text{GL}(\mathrm{\Sigma })`$ on $`G_n(\mathrm{\Sigma })`$ induces the action of $`\text{GL}(\mathrm{\Sigma })`$ on the tangent bundle $`TG_n(\mathrm{\Sigma })`$. It is easy to see that the only invariant of a tangent vector $`\xi `$ for this action is $`\text{rank}\overline{\xi }`$ (tangent vectors are just โdouble pointsโ or โpairs of infinitesimaly close pointsโ and number $`(n\text{rank}\overline{\xi })`$ is the infinitesimal version of the dimension of the intersection for a pair of points in the Grassmannian). Formula (4) implies:
$$\text{rank}[\xi _0,\xi _1]\mathrm{min}\{\text{rank}\overline{\xi }_0,\text{rank}\overline{\xi }_1\}.$$
### 11 Curves in the Grassmannian
Let $`tv(t)`$ be a germ at $`\overline{t}`$ of a smooth curve in the Grassmannian $`G_n(\mathrm{\Sigma })`$.
Definition. We say that the germ $`v()`$ is ample if $`v(t)v(\overline{t})=0t\overline{t}`$ and the operator-valued function $`t\pi _{v(t)v(\overline{t})}`$ has a pole at $`\overline{t}`$. We say that the germ $`v()`$ is regular if the function $`t\pi _{v(t)v(\overline{t})}`$ has a simple pole at $`\overline{t}`$. A smooth curve in $`G_n(\mathrm{\Sigma })`$ is called ample (regular) if all its germs are ample (regular).
Assume that $`\mathrm{\Sigma }=\{(x,y):x,y^n\}`$ is coordinatized in such a way that $`v(\overline{t})=\{(x,0):x^n\}`$. Then $`v(t)=\{(x,S_tx):x^n\}`$, where $`S(\overline{t})=0`$ and $`\pi _{v(t)v(\overline{t})}=\left(\begin{array}{cc}I& S_t^1\\ 0& 0\end{array}\right).`$ The germ $`v()`$ is ample if and only if the scalar function $`tdetS_t`$ has a finite order root at $`\overline{t}`$. The germ $`v()`$ is regular if and only if the matrix $`\dot{S}_{\overline{t}}`$ is not degenerate. More generally, the curve $`\tau \{(x,S_\tau x):x^n\}`$ is ample if and only if $`t`$ the function $`\tau det(S_\tau S_t)`$ has a finite order root at $`t`$. This curve is regular if and only if $`det\dot{S}_t0,t.`$ The intrinsic version of this coordinate characterization of regularity reads: the curve $`v()`$ is regular if and only if the map $`\overline{\dot{v}}(t)\text{Hom}(v(t),\mathrm{\Sigma }/v(t))`$ has rank $`n,t`$.
Coming back to the vector fields and their Jacobi curves (see Sec. 8) one can easily check that a vector field is ample (regular) if and only if its Jacobi curves are ample (regular).
Let $`v()`$ be an ample curve in $`G_n(\mathrm{\Sigma })`$. We consider the Laurent expansions at $`t`$ of the operator-valued function $`\tau \pi _{v(\tau )v(t)}`$,
$$\pi _{v(\tau )v(t)}=\underset{i=k_t}{\overset{m}{}}(\tau t)^i\pi _t^i+O(\tau t)^{m+1}.$$
Projectors of $`\mathrm{\Sigma }`$ on the subspace $`v(t)`$ form an affine subspace of $`\text{gl}(\mathrm{\Sigma })`$ (cf. Lemma II.1). This fact implies that $`\pi _t^0`$ is a projector of $`\mathrm{\Sigma }`$ on $`v(t)`$; in other words, $`\pi _t^0=\pi _{v^{}(t)v(t)}`$ for some $`v^{}(t)v(t)^{}`$. We thus obtain another curve $`tv^{}(t)`$ in $`G_n(\mathrm{\Sigma })`$, where $`\mathrm{\Sigma }=v(t)v^{}(t),t`$. The curve $`tv^{}(t)`$ is called the derivative curve of the ample curve $`v()`$.
The affine space $`\{\pi _{wv(t)}:wv(t)^{}\}`$ is a translation of the linear space $`๐(v(t))=\{๐ซ:\mathrm{\Sigma }v(t)๐ซ|_{v(t)}=0\}\text{gl}(\mathrm{\Sigma })\}`$ containing only nilpotent operators. It is easy to see that $`\pi _t^i๐(v(t))`$ for $`i0`$.
The derivative curve is not necessary ample. Moreover, it may be nonsmooth and even discontinuous.
###### Lemma II.3
If $`v()`$ is regular then $`v^{}()`$ is smooth.
Proof. Weโll find the coordinate representation of $`v^{}()`$. Let $`v(t)=\{(x,S_tx):x^n\}`$. Regularity of $`v()`$ is equivalent to the nondegeneracy of $`\dot{S}_t`$. We have:
$$\pi _{v(\tau )v(t)}=\left(\begin{array}{cc}S_{\tau t}^1S_\tau & S_{\tau t}^1\\ S_tS_{\tau t}^1S_\tau & S_tS_{\tau t}^1\end{array}\right),$$
where $`S_{\tau t}=S_\tau S_t`$. Then $`S_{\tau t}^1=(\tau t)^1\dot{S}_t^1\frac{1}{2}\dot{S}_t^1\ddot{S}_t\dot{S}_t^1+O(\tau t)`$ as $`\tau t`$ and
$$\pi _{v(\tau )v(t)}=(\tau t)^1\left(\begin{array}{cc}\dot{S}_t^1S_t& \dot{S}_t^1\\ S_t\dot{S}_t^1S_t& S_t\dot{S}_t^1\end{array}\right)+$$
$$\left(\begin{array}{cc}I\frac{1}{2}\dot{S}_t^1\ddot{S}_t\dot{S}_t^1S_t& \frac{1}{2}\dot{S}_t^1\ddot{S}_t\dot{S}_t^1\\ S_t\frac{1}{2}S_t\dot{S}_t^1\ddot{S}_t\dot{S}_t^1S_t& \frac{1}{2}S_t\dot{S}_t^1\ddot{S}_t\dot{S}_t^1\end{array}\right)+O(\tau t).$$
We set $`A_t=\frac{1}{2}\dot{S}_t^1\ddot{S}_t\dot{S}_t^1`$; then $`\pi _{v^{}(t)v(t)}=\left(\begin{array}{cc}I+A_tS_t& A_t\\ S_t+S_tA_tS_t& S_tA_t\end{array}\right)`$ is smooth with respect to $`t`$. Hence $`tv^{}(t)`$ is smooth. We obtain:
$$v^{}(t)=\{(A_ty,y+S_tA_ty):y^n\}.$$
$`(5)`$
### 12 The curvature
Definition. Let $`v`$ be an ample curve and $`v^{}`$ be the derivative curve of $`v`$. Assume that $`v^{}`$ is differentiable at $`t`$ and set $`R_v(t)=[\dot{v}^{}(t),\dot{v}(t)]`$. The operator $`R_v(t)gl(v(t))`$ is called the curvature of the curve $`v`$ at $`t`$.
If $`v`$ is a regular curve, then $`v^{}`$ is smooth, the curvature is well-defined and has a simple coordinate presentation. To find this presentation, weโll use formula (4) applied to $`\xi _0=\dot{v}^{}(t),\xi _1=\dot{v}(t)`$. As before, we assume that $`v(t)=\{(x,S_tx):x^n\}`$; in particular, $`v(t)`$ is transversal to the subspace $`\{(0,y):y^n\}`$. In order to apply (4) we need an extra assumption on the coordinatization of $`\mathrm{\Sigma }`$: the subspace $`v^{}(t)`$ has to be transversal to $`\{(0,y):y^n\}`$ for given $`t`$. The last property is equivalent to the nondegeneracy of the matrix $`A_t`$ (see (6)). It is important to note that the final expression for $`R_v(t)`$ as a differential operator of $`S`$ must be valid without this extra assumption since the definition of $`R_v(t)`$ is intrinsic! Now we compute: $`v^{}(t)=\{(x,(A_t^1+S_t)x):x^n\},R_v(t)=[\dot{v}^{}(t),\dot{v}(t)]=A_t\frac{d}{dt}(A_t^1+S_t)A_t\dot{S}_t=(A_t\dot{S}_t)^2\dot{A}_t\dot{S}_t=\frac{1}{4}(\dot{S}_t^1\ddot{S}_t)^2\dot{A}_t\dot{S}_t.`$ We also have $`\dot{A}\dot{S}=\frac{1}{2}\frac{d}{dt}(\dot{S}^1\ddot{S}\dot{S}^1)\dot{S}=(\dot{S}^1)^2\frac{1}{2}\dot{S}^1\stackrel{\mathrm{}}{S}`$. Finally,
$$R_v(t)=\frac{1}{2}\dot{S}_t^1\underset{t}{\overset{\mathrm{}}{S}}\frac{3}{4}(\dot{S}_t^1\ddot{S}_t)^2=\frac{d}{dt}\left((2\dot{S}_t)^1\ddot{S}_t\right)\left((2\dot{S}_t)^1\ddot{S}_t\right)^2,$$
$`(6)`$
the matrix version of the Schwartzian derivative.
Curvature operator is a fundamental invariant of the curve in the Grassmannian. One more intrinsic construction of this operator, without using the derivative curve, is provided by the following
###### Proposition II.1
Let $`v`$ be a regular curve in $`G_n(\mathrm{\Sigma })`$. Then
$$[\dot{v}(\tau ),\dot{v}(t)]=(\tau t)^2\text{id}+\frac{1}{3}R_v(t)+O(\tau t)$$
as $`\tau t`$.
Proof. It is enough to check the identity in some coordinates. Given $`t`$ we may assume that
$$v(t)=\{(x,0):x^n\},v^{}(t)=\{(0,y):y^n\}.$$
Let $`v(\tau )=\{(x,S_\tau x:x^n\}`$, then $`S_t=\ddot{S}_t=0`$ (see (5)). Moreover, we may assume that the bases of the subspaces $`v(t)`$ and $`v^{}(t)`$ are coordinated in such a way that $`\dot{S}_t=I`$. Then $`R_v(t)=\frac{1}{2}\underset{t}{\overset{\mathrm{}}{S}}`$ (see (6)). On the other hand, formula (4) for the infinitesimal cross-ratio implies:
$$[\dot{v}(\tau ),\dot{v}(t)]=S_\tau ^1\dot{S}_\tau S_\tau ^1=\frac{d}{d\tau }(S_\tau ^1)=$$
$$\frac{d}{d\tau }((\tau t)I+\frac{(\tau t)^3}{6}\underset{t}{\overset{\mathrm{}}{S}})^1+O(\tau t)=$$
$$\frac{d}{d\tau }((\tau t)^1I\frac{(\tau t)}{6}\underset{t}{\overset{\mathrm{}}{S}})+O(\tau t)=(\tau t)^2I+\frac{1}{6}\underset{t}{\overset{\mathrm{}}{S}}+O(\tau t).$$
$`\mathrm{}`$
Curvature operator is an invariant of the curves in $`G_n(\mathrm{\Sigma })`$ with fixed parametrizations. Asymptotic presentation obtained in Proposition II.1 implies a nice chain rule for the curvature of the reparametrized curves.
Let $`\phi :`$ be a regular change of variables, i.e. $`\dot{\phi }0,t`$. The standard imbedding $`^1=G_1(^2)`$ makes $`\phi `$ a regular curve in $`G_1(^2)`$. As we know (see (6)), the curvature of this curve is the Schwartzian of $`\phi `$:
$$R_\phi (t)=\frac{\stackrel{\mathrm{}}{\phi }(t)}{2\dot{\varphi }(t)}\frac{3}{4}\left(\frac{\ddot{\phi }(t)}{\dot{\phi }(t)}\right)^2.$$
We set $`v_\phi (t)=v(\phi (t))`$ for any curve $`v`$ in $`G_n(\mathrm{\Sigma })`$.
###### Proposition II.2
Let $`v`$ be a regular curve in $`G_n(\mathrm{\Sigma })`$ and $`\phi :`$ be a regular change of variables. Then
$$R_{v_\phi }(t)=\dot{\phi }^2(t)R_v(\phi (t))+R_\phi (t).$$
$`(7)`$
Proof. We have
$$[\dot{v}_\phi (\tau ),\dot{v}_\phi (t)]=(\tau t)^2\text{id}+\frac{1}{3}R_{v_\phi }(t)+O(\tau t).$$
On the other hand,
$$[\dot{v}_\phi (\tau ),\dot{v}_\phi (t)]=[\dot{\phi }(\tau )\dot{v}(\phi (\tau )),\dot{\phi }(t)\dot{v}(\phi (t))]=\dot{\phi }(\tau )\dot{\phi }(t)[\dot{v}(\phi (\tau )),\dot{v}(\phi (t))]=$$
$$\dot{\phi }(\tau )\dot{\phi }(t)\left((\phi (\tau )\phi (t))^2id+\frac{1}{3}R_v(\phi (t))+O(\tau t)\right)=$$
$$\frac{\dot{\phi }(\tau )\dot{\phi }(t)}{((\phi (\tau )\phi (t))^2}\text{id}+\frac{\dot{\phi }^2(t)}{3}R_v(\phi (t))+O(\tau t).$$
We treat $`\phi `$ as a curve in $`^1=G_1(^2)`$. Then $`[\dot{\phi }(\tau ),\dot{\phi }(t)]=\frac{\dot{\phi }(\tau )\dot{\phi }(t)}{(\phi (\tau )\phi (t))^2}`$, see (4). The one-dimensional version of Proposition II.1 reads:
$$[\dot{\phi }(\tau ),\dot{\phi }(t)]=(t\tau )^2+\frac{1}{3}R_\phi (t)+O(\tau t).$$
Finally,
$$[\dot{v}_\phi (\tau ),\dot{v}_\phi (t)]=(t\tau )^2+\frac{1}{3}(R_\phi (t)+\dot{\phi }^2(t)R_v(\phi (t)))+O(\tau t).\mathrm{}$$
The following identity is an immediate corollary of Proposition II.2:
$$\left(R_{v_\phi }\frac{1}{n}(\text{tr}R_{v_\phi })\text{id}\right)(t)=\dot{\phi }^2(t)\left(R_v\frac{1}{n}(\text{tr}R_v)\text{id}\right)(\phi (t)).$$
$`(8)`$
Definition. An ample curve $`v`$ is called flat if $`R_v(t)0`$.
It follows from Proposition II.1 that any small enough piece of a regular curve can be made flat by a reparametrization if and only if the curvature of the curve is a scalar operator, i.e. $`R_v(t)=\frac{1}{n}(\text{tr}R_v(t))\text{id}`$. In the case of a nonscalar curvature, one can use equality (8) to define a distinguished parametrization of the curve and then derive invariants which do not depend on the parametrization.
Remark. In this paper we are mainly focused on the regular curves. See paper for the version of the chain rule which is valid for any ample curve and for basic invariants of unparametrized ample curves.
### 13 Structural equations
Assume that $`v`$ and $`w`$ are two smooth curves in $`G_n(\mathrm{\Sigma })`$ such that $`v(t)w(t)=0,t`$.
###### Lemma II.4
For any $`t`$ and any $`ev(t)`$ there exists a unique $`f_ew(t)`$ with the following property: $``$ a smooth curve $`e_\tau v(\tau ),e_t=e`$, such that $`\frac{d}{d\tau }e_\tau |_{\tau =t}=f_e`$. Moreover, the mapping $`\mathrm{\Phi }_t^{vw}:ef_t`$ is linear and for any $`e_0v(0)`$ there exists a unique smooth curve $`e(t)v(t)`$ such that $`e(0)=e_0`$ and
$$\dot{e}(t)=\mathrm{\Phi }_t^{vw}e(t),t.$$
$`(9)`$
Proof. First we take any curve $`\widehat{e}_\tau v(\tau )`$ such that $`e_t=e`$. Then $`\widehat{e}_\tau =a_\tau +b_\tau `$ where $`a_\tau v(t),b_\tau w(t)`$. We take $`x_\tau v(\tau )`$ such that $`x_t=\dot{a}_t`$ and set $`e_\tau =\widehat{e}_\tau +(t\tau )x_\tau `$. Then $`\dot{e}_t=\dot{b}_t`$ and we put $`f_e=\dot{b}_t`$.
Let us prove that $`\dot{b}_t`$ depends only on $`e`$ and not on the choice of $`e_\tau `$. Computing the difference of two admissible $`e_\tau `$ we reduce the lemma to the following statement: if $`z(\tau )v(\tau ),\tau `$ and $`z(t)=0`$, then $`\dot{z}(t)v(t)`$.
To prove the last statement we take smooth vector-functions $`e_\tau ^iv(\tau ),i=1,\mathrm{},n`$ such that $`v(\tau )=span\{e_\tau ^1,\mathrm{},e_\tau ^n\}`$. Then $`z(\tau )=\underset{i=1}{\overset{n}{}}\alpha _i(\tau )e_\tau ^i,\alpha _i(t)=0`$. Hence $`\dot{z}(t)=\underset{i=1}{\overset{n}{}}\dot{\alpha }_i(t)e_t^iv_t.`$
Linearity of the map $`\mathrm{\Phi }_t^{vw}`$ follows from the uniqueness of $`f_e`$. Indeed, if $`f_{e^i}=\frac{d}{d\tau }e_\tau ^i|_{\tau =t}`$, then $`\frac{d}{d\tau }(\alpha _1e_\tau ^1+\alpha _2e_\tau ^2)|_{\tau =t}=\alpha _1f_{e^1}+\alpha _2f_{e^2}`$; hence $`\alpha _1f_{e^1}+\alpha _2f_{e^2}=f_{\alpha _1e^1+\alpha _2e^2},e^iv(t),\alpha _i,i=1,2`$.
Now consider the smooth submanifold $`V=\{(t,e):t,ev(t)\}`$ of $`\times \mathrm{\Sigma }`$. We have $`(1,\mathrm{\Phi }_t^{vw}e)T_{(t,e)}V`$ since $`(1,\mathrm{\Phi }_t^{vw}e)`$ is the velocity of a curve $`\tau (\tau ,e_\tau )`$ in $`V`$. So $`(t,e)(1,\mathrm{\Phi }_t^{vw}e),(t,e)V`$ is a smooth vector field on $`V`$. The curve $`e(t)v(t)`$ satisfies (9) if and only if $`(t,e(t))`$ is a trajectory of this vector field. Now the standard existence and uniqueness theorem for ordinary differential equations provides the existence of a unique solution to the Cauchy problem for small enough $`t`$ while the linearity of the equation guarantees that the solution is defined for all $`t.\mathrm{}`$
It follows from the proof of the lemma that $`\mathrm{\Phi }_t^{vw}e=\pi _{v(t)w(t)}\dot{e}_\tau |_{\tau =t}`$ for any $`e_\tau v(\tau )`$ such that $`v_t=e`$. Let $`v(t)=\{(x,S_{vt}x):x^n\},w(t)=\{(x,S_{wt}x):x^n\}`$; the matrix presentation of $`\mathrm{\Phi }_t^{vw}`$ in coordinates $`x`$ is $`(S_{wt}S_{vt})^1\dot{S}_{vt}`$. Linear mappings $`\mathrm{\Phi }_t^{vw}`$ and $`\mathrm{\Phi }_t^{wv}`$ provide a factorization of the infinitesimal cross-ratio $`[\dot{w}(t),\dot{v}(t)]`$. Indeed, equality (4) implies:
$$[\dot{w}(t),\dot{v}(t)]=\mathrm{\Phi }_t^{wv}\mathrm{\Phi }_t^{vw}.$$
$`(10)`$
Equality (9) implies one more useful presentation of the infinitesimal cross-ratio: if $`e(t)`$ satisfies (9), then
$$[\dot{w}(t),\dot{v}(t)]e(t)=\mathrm{\Phi }_t^{wv}\mathrm{\Phi }_t^{vw}e(t)=\mathrm{\Phi }_t^{wv}\dot{e}(t)=\pi _{w(t)v(t)}\ddot{e}(t).$$
$`(11)`$
Now let $`w`$ be the derivative curve of $`v`$, $`w(t)=v^{}(t)`$. It happens that $`\ddot{e}(t)v(t)`$ in this case and (11) is reduced to the structural equation:
$$\ddot{e}(t)=[\dot{v}^{}(t),\dot{v}(t)]e(t)=R_v(t)e(t),$$
where $`R_v(t)`$ is the curvature operator. More precisely, we have the following
###### Proposition II.3
Assume that $`v`$ is a regular curve in $`G_n(\mathrm{\Sigma })`$, $`v^{}`$ is its derivative curve, and $`e()`$ is a smooth curve in $`\mathrm{\Sigma }`$ such that $`e(t)v(t),t`$. Then $`\dot{e}(t)v^{}(t)`$ if and only if $`\ddot{e}(t)v(t).`$
Proof. Given $`t`$, we take coordinates in such a way that $`v(t)=\{(x,0):x^n\},v^{}(t)=\{(0,y):y^n\}`$. Then $`v(\tau )=\{(x,S_\tau x):x^n\}`$ for $`\tau `$ close enough to $`t`$, where $`S_t=\ddot{S}_t=0`$ (see (5)).
Let $`e(\tau )=\{(x(\tau ),S_\tau x(\tau ))\}`$. The inclusion $`\dot{e}(t)v^{}(t)`$ is equivalent to the equality $`\dot{x}(t)=0`$. Further,
$$\ddot{e}(t)=\{\ddot{x}(t),\ddot{S}_tx(t)+2\dot{S}_t\dot{x}(t)+S_t\ddot{x}(t)\}=\{\ddot{x}(t),2\dot{S}\dot{x}\}v(t).$$
Regularity of $`v`$ implies the nondegeneracy of $`\dot{S}(t)`$. Hence $`\ddot{e}(t)v(t)`$ if and only if $`\dot{x}(t)=0.\mathrm{}`$
Now equality (11) implies
###### Corollary II.1
If $`\dot{e}(t)=\mathrm{\Phi }_t^{vv^{}}e(t)`$, then $`\ddot{e}(t)+R_v(t)e(t)=0`$.
Let us consider invertible linear mappings $`V_t:v(0)v(t)`$ defined by the relations $`V_te(0)=e(t),\dot{e}(\tau )=\mathrm{\Phi }_\tau ^{vv^{}}e(\tau ),0\tau t`$. It follows from the structural equation that the curve $`v`$ is uniquely reconstructed from $`\dot{v}(0)`$ and the curve $`tV_t^1R_V(t)`$ in $`\text{gl}(v(0))`$. Moreover, let $`v_0G_n(\mathrm{\Sigma })`$ and $`\xi T_{v_0}G_n(\mathrm{\Sigma })`$, where the map $`\overline{\xi }\text{Hom}(v_0,\mathrm{\Sigma }/v_0)`$ has rank $`n`$; then for any smooth curve $`tA(t)`$ in $`\text{gl}(v_0)`$ there exists a unique regular curve $`v`$ such that $`\dot{v}(0)=\xi `$ and $`V_t^1R_v(t)V_t=A(t)`$. Indeed, let $`e_i(0),i=1,\mathrm{},n`$, be a basis of $`v_0`$ and $`A(t)e_i(0)=\underset{j=1}{\overset{n}{}}a_{ij}(t)e_j(0)`$. Then $`v(t)=span\{e_1(t),\mathrm{},e_n(t)\}`$, where
$$\ddot{e}_i(\tau )+\underset{j=1}{\overset{n}{}}a_{ij}(\tau )e_j(\tau )=0,0\tau t,$$
$`(12)`$
are uniquely defined by fixing the $`\dot{v}(0)`$.
The obtained classification of regular curves in terms of the curvature is particularly simple in the case of a scalar curvature operators $`R_v(t)=\rho (t)\text{id}`$. Indeed, we have $`A(t)=V_t^1R_v(t)V_t=\rho (t)\text{id}`$ and system (12) is reduced to $`n`$ copies of the Hill equation $`\ddot{e}(\tau )+\rho (\tau )e(\tau )=0`$.
Recall that all $`\xi TG_n(\mathrm{\Sigma })`$ such that $`\text{rank}\overline{\xi }=n`$ are equivalent under the action of $`\text{GL}(\mathrm{\Sigma })`$ on $`TG_n(\mathrm{\Sigma })`$ induced by the standard action on the Grassmannian $`G_n(\mathrm{\Sigma })`$. We thus obtain
###### Corollary II.2
For any smooth scalar function $`\rho (t)`$ there exists a unique, up to the action of $`\text{GL}(\mathrm{\Sigma })`$, regular curve $`v`$ in $`G_n(\mathrm{\Sigma })`$ such that $`R_v(t)=\rho (t)\text{id}`$.
Another important special class is that of symmetric curves.
Definition. A regular curve $`v`$ is called symmetric if $`V_tR_v(t)=R_v(t)V_t,t`$.
In other words, $`v`$ is symmetric if and only the curve $`A(t)=V_t^1R_v(t)V_t`$ in $`\text{gl}(v(0))`$ is constant and coincides with $`R_v(0)`$. The structural equation implies
###### Corollary II.3
For any $`n\times n`$-matrix $`A_0`$, there exists a unique, up to the action of $`\text{GL}(\mathrm{\Sigma })`$, symmetric curve $`v`$ such that $`R_v(t)`$ is similar to $`A_0`$.
The derivative curve $`v^{}`$ of a regular curve $`v`$ is not necessary regular. The formula $`R_v(t)=\mathrm{\Phi }_t^{v^{}v}\mathrm{\Phi }_t^{vv^{}}`$ implies that $`v^{}`$ is regular if and only if the curvature operator $`R_v(t)`$ is nondegenerate for any $`t`$. Then we may compute the second derivative curve $`v^{}=(v^{})^{}`$.
###### Proposition II.4
A regular curve $`v`$ with nondegenerate curvature operators is symmetric if and only if $`v^{}=v`$.
Proof. Let us consider system (12). We are going to apply Proposition II.3 to the curve $`v^{}`$ (instead of $`v`$) and the vectors $`\dot{e}_i(t)v^{}(t)`$. According to Proposition II.3, $`v^{}=v`$ if and only if $`\frac{d^2}{dt^2}\dot{e}_i(t)v^{}(t)`$. Differentiating (12) we obtain that $`v^{}=v`$ if and only if the functions $`\alpha _{ij}(t)`$ are constant. The last property is none other than a characterization of symmetric curves.$`\mathrm{}`$
### 14 Canonical connection
Now we apply the developed theory of curves in the Grassmannian to the Jacobi curves $`J_z(t)`$ (see Sec. 8).
###### Proposition II.5
All Jacobi curves $`J_z(),zN`$, associated to the given vector field $`\zeta `$ are regular (ample) if and only if the field $`\zeta `$ is regular (ample).
Proof. The definition of the regular (ample) field is actually the specification of the definition of the regular (ample) germ of the curve in the Grassmannian: general definition is applied to the germs at $`t=0`$ of the curves $`tJ_z(t)`$. What remains is to demonstrate that other germs of these curves are regular (ample) as soon as the germs at 0 are. The latter fact follows from the identity
$$J_z(t+\tau )=e_{}^{t\zeta }J_{e^{t\zeta }(z)}(\tau )$$
$`(13)`$
(which, in turn, is an immediate corollary of the identity $`e_{}^{(t+\tau )\zeta }=e_{}^{t\zeta }e_{}^{\tau \zeta }`$). Indeed, (13) implies that the germ of $`J_z()`$ at $`t`$ is the image of the germ of $`J_{e^{t\zeta }(\tau )}()`$ at 0 under the fixed linear transformation $`e_{}^{t\zeta }:T_{e^{t\zeta }(z)}NT_zN`$. The properties of the germs to be regular or ample survive linear transformations since they are intrinsic properties. $`\mathrm{}`$
Let $`\zeta `$ be an ample field. Then the derivative curves $`J_z^{}(t)`$ are well-defined. Moreover, identity (13) and the fact that the construction of the derivative curve is intrinsic imply:
$$J_z^{}(t)=e_{}^{t\zeta }J_{e^{t\zeta }(z)}^{}(0).$$
$`(14)`$
The value at 0 of the derivative curve provides the splitting $`T_zM=J_z(0)J_z^{}(0)`$, where the first summand is the tangent space to the fiber, $`J_z(0)=T_zE_z`$.
Now assume that $`J_z^{}(t)`$ smoothly depends on $`z`$; this assumption is automatically fulfilled in the case of a regular $`\zeta `$, where we have the explicit coordinate presentation for $`J_z^{}(t)`$. Then the subspaces $`J_z^{}(0)T_zN,zN,`$ form a smooth vector distribution, which is the direct complement to the vertical distribution $`=\{T_zE_z:zN\}`$. Direct complements to the vertical distribution are called Ehresmann connections (or just nonlinear connections, even if linear connections are their special cases). The Ehresmann connection $`_\zeta =\{J_z^{}(0):zN\}`$ is called the canonical connection associated with $`\zeta `$ and the correspondent splitting $`TN=_\zeta `$ is called the canonical splitting. Our nearest goal is to give a simple intrinsic characterization of $`_\zeta `$ which does not require the integration of the equation $`\dot{z}=\zeta (z)`$ and is suitable for calculations not only in local coordinates but also in moving frames.
Let $`=\{F_zT_zN:zN\}`$ be an Ehresmann connection. Given a vector field $`\xi `$ on $`E`$ we denote $`\xi _{ver}(z)=\pi _{F_zJ_z(0)}\xi ,\xi _{hor}(z)=\pi _{J_z(0)F_z}\xi `$, the โverticalโ and the โhorizontalโ parts of $`\xi (z)`$. Then $`\xi =\xi _{ver}+\xi _{hor}`$, where $`\xi _{ver}`$ is a section of the distribution $``$ and $`\xi _{hor}`$ is a section of the distribution $``$. In general, sections of $``$ are called vertical fields and sections of $``$ are called horizontal fields.
###### Proposition II.6
Assume that $`\zeta `$ is a regular field. Then $`=_\zeta `$ if and only if the equality
$$[\zeta ,[\zeta ,\nu ]]_{hor}=2[\zeta ,[\zeta ,\nu ]_{ver}]_{hor}$$
$`(15)`$
holds for any vertical vector field $`\nu `$. Here $`[,]`$ is Lie bracket of vector fields.
Proof. The deduction of identity (15) is based on the following classical expression:
$$\frac{d}{dt}e_{}^{t\zeta }\xi =e_{}^{t\zeta }[\zeta ,\xi ],$$
$`(16)`$
for any vector field $`\xi `$.
Given $`zN`$, we take coordinates in $`T_zN`$ in such a way that $`T_zN=\{(x,y):x,y^n\}`$, where $`J_z(0)=\{(x,0):x^n\},J_z^{}(0)=\{(0,y):y^n\}`$. Let $`J_z(t)=\{(x,S_tx):x^n\}`$, then $`S_0=\ddot{S}_0=0`$ and $`det\dot{S}_00`$ due to the regularity of the Jacobi curve $`J_z`$.
Let $`\nu `$ be a vertical vector field, $`\nu (z)=(x_0,0)`$ and $`\left(e_{}^{t\zeta }\nu \right)(z)=(x_t,y_t)`$. Then $`(x_t,0)=\left(e_{}^{t\zeta }\nu \right)_{ver}(z),(0,y_t)=\left(e_{}^{t\zeta }\nu \right)_{hor}(z)`$. Moreover, $`y_t=S_tx_t`$ since $`\left(e_{}^{t\zeta }\nu \right)(z)J_z(t)`$. Differentiating the identity $`y_t=S_tx_t`$ we obtain: $`\dot{y}_t=\dot{S}_tx_t+S_t\dot{x}_t.`$ In particular, $`\dot{y}_0=\dot{S}_0x_0`$. It follows from (16) that $`(\dot{x}_0,0)=[\zeta ,\nu ]_{ver},(0,\dot{y}_0)=[\zeta ,\nu ]_{hor}`$. Hence $`(0,\dot{S}_0x_0)=[\zeta ,\nu ]_{hor}(z)`$, where, I recall, $`\nu `$ is any vertical field. Now we differentiate once more and evaluate the derivative at 0:
$$\ddot{y}_0=\ddot{S}_0x_0+2\dot{S}_0\dot{x}_0+S_0\ddot{x}_0=2\dot{S}_0\dot{x}_0.$$
$`(17)`$
The Lie bracket presentations of the left and right hand sides of (17) are: $`(0,\ddot{y}_0)=[\zeta ,[\zeta ,\nu ]]_{hor},(0,\dot{S}_0\dot{x}_0)=[\zeta ,[\zeta ,\nu ]_{ver}]_{hor}`$. Hence (17) implies identity (15).
Assume now that $`\{(0,y):y^n\}J_z^{}(0)`$; then $`\ddot{S}_0x_00`$ for some $`x_0`$. Hence $`\ddot{y}_02\dot{S}_0\dot{x}_0`$ and equality (15) is violated. $`\mathrm{}`$
Inequality (15) can be equivalently written in the following form that is often more convenient for the computations:
$$\pi _{}[\zeta ,[\zeta ,\nu ]](z)=2\pi _{}[\zeta ,[\zeta ,\nu ]_{ver}](z),zN.$$
$`(18)`$
Let $`R_{J_z}(t)\text{gl}(J_z(t))`$ be the curvature of the Jacobi curve $`J_z(t)`$. Identity (13) and the fact that construction of the Jacobi curve is intrinsic imply that
$$R_{J_z}(t)=e_{}^{t\zeta }R_{J_{e^{t\zeta }(z)}}(0)e_{}^{t\zeta }|_{J_z(t)}.$$
Recall that $`J_z(0)=T_zE_z`$; the operator $`R_{J_z}(0)\text{gl}(T_zE_z)`$ is called the curvature operator of the field $`\zeta `$ at $`z`$. We introduce the notation: $`R_\zeta (z)\stackrel{def}{=}R_{J_z}(0)`$; then $`R_\zeta =\left\{R_\zeta (z)\right\}_{zE}`$ is an automorphism of the โverticalโ vector bundle $`\left\{T_zE_z\right\}_{zM}`$.
###### Proposition II.7
Assume that $`\zeta `$ is an ample vector field and $`J_z^{}(0)`$ is smooth with respect to $`z`$. Let $`TN=_\zeta `$ be the canonical splitting. Then
$$R_\zeta \nu =[\zeta ,[\zeta ,\nu ]_{hor}]_{ver}$$
$`(19)`$
for any vertical field $`\nu `$.
Proof. Recall that $`R_{J_z}(0)=[\dot{J}_z^{}(0),\dot{J}_z(0)]`$, where $`[,]`$ is the infinitesimal crossโratio (not the Lie bracket!). The presentation (10) of the infinitesimal crossโratio implies:
$$R_\zeta (z)=R_{J_z}(0)=\mathrm{\Phi }_0^{J_z^{}J_z}\mathrm{\Phi }_0^{J_zJ_z^{}},$$
where $`\mathrm{\Phi }_0^{vw}e=\pi _{v(0)w(0)}\dot{e}_0`$ for any smooth curve $`e_\tau v(\tau )`$ such that $`e_0=e`$. Equalities (14) and (16) imply: $`\mathrm{\Phi }_0^{J_zJ_z^{}}\nu (z)=[\zeta ,\nu ]_{ver}(z),zM.`$ Similarly, $`\mathrm{\Phi }_0^{J_z^{}J_z}\mu (z)=[\zeta ,\mu ]_{hor}(z)`$ for any horizontal field $`\mu `$ and any $`zM`$. Finally,
$$R_\zeta (z)\nu (z)=\mathrm{\Phi }_0^{J_z^{}J_z}\mathrm{\Phi }_0^{J_zJ_z^{}}=[\zeta ,[\zeta ,\nu ]_{hor}]_{ver}(z).$$
$`\mathrm{}`$
### 15 Coordinate presentation
We fix local coordinates acting in the domain $`๐ชN`$, which turn the foliation into the Cartesian product of vector spaces: $`๐ช\{(x,y):x,y^n\}`$, $`\pi :(x,y)y`$. Then vector field $`\zeta `$ takes the form $`\zeta =\underset{i=1}{\overset{n}{}}\left(a^i\frac{}{x_i}+b^i\frac{}{y_i}\right)`$, where $`a^i,b^i`$ are smooth functions on $`^n\times ^n`$. Below we use abridged notations: $`\frac{}{x_i}=_{x_i},\frac{\phi }{x_i}=\phi _{x_i}`$ etc. We also use the standard summation agreement for repeating indices.
Recall the coordinate characterization of the regularity property for the vector field $`\zeta `$. Intrinsic definition of regular vector fields is done in Section 8; it is based on the mapping $`\mathrm{\Pi }_z`$ whose coordinate presentation is: $`\mathrm{\Pi }_{(x,y)}:x(b^1(x,y),\mathrm{},b^n(x,y))^{}`$. Field $`\zeta `$ is regular if and only if $`\mathrm{\Pi }_y`$ are submersions; in other words, if and only if $`\left(b_{x_j}^i\right)_{i,j=1}^n`$ is a non degenerate matrix.
Vector fields $`_{x_i},i=1,\mathrm{},n`$, provide a basis of the space of vertical fields. As soon as coordinates are fixed, any Ehresmann connection finds a unique basis of the form:
$$\left(_{y_i}\right)_{hor}=_{y_i}+c_i^j_{x_j},$$
where $`c_i^j,i,j=1,\mathrm{},n`$, are smooth functions on $`^n\times ^n`$. To characterize a connection in coordinates thus means to find functions $`c_i^j`$. In the case of the canonical connection of a regular vector field, the functions $`c_i^j`$ can be easily recovered from identity (18) applied to $`\nu =_{x_i},i=1,\mathrm{},n`$. Weโll do it explicitly for two important classes of vector fields: second order ordinary differential equations and Hamiltonian systems.
A second order ordinary differential equation
$$\dot{y}=x,\dot{x}=f(x,y)$$
$`(20)`$
there corresponds to the vector field $`\zeta =f^i_{x_i}+x_i_{y_i}`$, where $`f=(f_1,\mathrm{},f_n)^{}`$. Let $`\nu =_{x_i}`$; then
$$[\zeta ,\nu ]=_{y_i}f_{x_i}^j_{x_j},[\zeta ,\nu ]_{ver}=(c_i^jf_{x_i}^j)_{x_j},$$
$$\pi _{}[\zeta ,[\zeta ,\nu ]]=f_{x_i}^j_{y_j},\pi _{}[\zeta ,[\zeta ,\nu ]_{ver}]=(f_{x_i}^jc_i^j)_{y_j}.$$
Hence, in virtue of equality (18) we obtain that $`c_i^j=\frac{1}{2}f_{x_i}^j`$ for the canonical connection associated with the second order differential equation (20).
Now consider a Hamiltonian vector field $`\zeta =h_{y_i}_{x_i}+h_{x_i}_{y_i}`$, where $`h`$ is a smooth function on $`^n\times ^n`$ (a Hamiltonian). The field $`\zeta `$ is regular if and only if the matrix $`h_{xx}=\left(h_{x_ix_j}\right)_{i,j=1}^n`$ is non degenerate. We are going to characterize the canonical connection associated with $`\zeta `$. Let $`C=\left(c_i^j\right)_{i,j=1}^n`$; the straightforward computation similar to the computation made for the second order ordinary differential equation gives the following presentation for the matrix $`C`$:
$$2\left(h_{xx}Ch_{xx}\right)_{ij}=h_{x_k}h_{x_ix_jy_k}h_{y_k}h_{x_ix_jx_k}h_{x_iy_k}h_{x_kx_j}h_{x_ix_k}h_{y_kx_j}$$
or, in the matrix form:
$$2h_{xx}Ch_{xx}=\{h,h_{xx}\}h_{xy}h_{xx}h_{xx}h_{yx},$$
where $`\{h,h_{xx}\}`$ is the Poisson bracket: $`\{h,h_{xx}\}_{ij}=\{h,h_{x_ix_j}\}=h_{x_k}h_{x_ix_jy_k}h_{y_k}h_{x_ix_jx_k}`$.
Note that matrix $`C`$ is symmetric in the Hamiltonian case (indeed, $`h_{xx}h_{yx}=(h_{xy}h_{xx})^{}`$). This is not occasional and is actually guaranteed by the fact that Hamiltonian flows preserve symplectic form $`dx_idy_i`$. See Section 17 for the symplectic version of the developed theory.
As soon as we found the canonical connection, formula (19) gives us the presentation of the curvature operator although the explicit coordinate expression can be bulky. Let us specify the vector field more. In the case of the Hamiltonian of a natural mechanical system, $`h(x,y)=\frac{1}{2}|x|^2+U(y)`$, the canonical connection is trivial: $`c_i^j=0`$; the matrix of the curvature operator is just $`U_{yy}`$.
Hamiltonian vector field associated to the Hamiltonian$`h(x,y)=g^{ij}(y)x_ix_j`$ with a non degenerate symmetric matrix $`\left(g^{ij}\right)_{i,j=1}^n`$ generates a (pseudo-)Riemannian geodesic flow. Canonical connection in this case is classical Levi Civita connection and the curvature operator is Ricci operator of (pseudo-)Riemannian geometry (see \[4, Sec. 5\] for details). Finally, Hamiltonian $`h(x,y)=g^{ij}(y)x_ix_j+U(y)`$ has the same connection as Hamiltonion $`h(x,y)=g^{ij}(y)x_ix_j`$ while its curvature operator is sum of Ricci operator and second covariant derivative of $`U`$.
### 16 Affine foliations
Let $`\left[\right]`$ be the sheaf of germs of sections of the distribution $`=\{T_zE_z:zN\}`$ equipped with the Lie bracket operation. Then $`\left[\right]_z`$ is just the Lie algebra of germs at $`zM`$ of vertical vector fields. Affine structure on the foliation $`E`$ is a sub-sheaf $`\left[\right]^a\left[\right]`$ such that $`\left[\right]_z^a`$ is an Abelian sub-algebra of $`\left[\right]_z`$ and $`\{\varsigma (z):\varsigma \left[\right]_z^a\}=T_zE_z,zN`$. A foliation with a fixed affine structure is called the affine foliation.
The notion of the affine foliation generalizes one of the vector bundle. In the case of the vector bundle, the sheaf $`\left[\right]^a`$ is formed by the germs of vertical vector fields whose restrictions to the fibers are constant (i.e. translation invariant) vector fields on the fibers. In the next section we will describe an important class of affine foliations which is not reduced to the vector bundles.
###### Lemma II.5
Let $``$ be an affine foliation, $`\varsigma \left[\right]_z^a`$ and $`\varsigma (z)=0`$. Then $`\varsigma |_{E_z}=0`$.
Proof. Let $`\varsigma _1,\mathrm{},\varsigma _n\left[\right]_z^a`$ be such that $`\varsigma _1(z),\mathrm{},\varsigma _n(z)`$ form a basis of $`T_zE_z`$. Then $`\varsigma =b_1\varsigma _1+\mathrm{}+b_n\varsigma _n`$, where $`b_i`$ are germs of smooth functions vanishing at $`z`$. Commutativity of $`\left[\right]_z^a`$ implies: $`0=[\varsigma _i,\varsigma ]=(\varsigma _ib_1)\varsigma _1+\mathrm{}+(\varsigma _ib_n)\varsigma _n`$. Hence functions $`b_i|_{E_z}`$ are constants, i.e. $`b_i|_{E_z}=0,i=1,\mathrm{},n.\mathrm{}`$
Lemma II.5 implies that $`\varsigma \left[\right]_z^a`$ is uniquely reconstructed from $`\varsigma (z)`$. This property permits to define the vertical derivative of any vertical vector field $`\nu `$ on $`M`$. Namely, $`vT_zE_z`$ we set
$$D_v\nu =[\varsigma ,\nu ](z),\text{where}\varsigma \left[\right]_z^a,\varsigma (z)=v.$$
Suppose $`\zeta `$ is a regular vector field on the manifold $`N`$ endowed with the affine $`n`$-foliation. The canonical Ehresmann connection $`_\zeta `$ together with the vertical derivative allow to define a canonical linear connection $``$ on the vector bundle $``$. Sections of the vector bundle $``$ are just vertical vector fields. We set
$$_\xi \nu =[\xi ,\nu ]_{ver}+D_\nu (\xi _{ver}),$$
where $`\xi `$ is any vector field on $`N`$ and $`\nu `$ is a vertical vector field. It is easy to see that $``$ satisfies axioms of a linear connection. The only non evident one is: $`_{b\xi }\nu =b_\xi \nu `$ for any smooth function $`b`$. Let $`zN,\varsigma \left[\right]_z^a`$, and $`\varsigma (z)=\nu (z)`$. We have
$$_{b\xi }\nu =[b\xi ,\nu ]_{ver}+[\varsigma ,b\xi _{ver}]=$$
$$b\left([\xi ,\nu ]_{ver}+[\varsigma ,\xi _{ver}]\right)(\nu b)\xi _{ver}+(\varsigma b)\xi _{ver}.$$
Hence
$$(_{b\xi }\nu )(z)=b(z)\left([\xi ,\nu ]_{ver}(z)+[\varsigma ,\xi _{ver}](z)\right)=(b_\xi \nu )(z).$$
Linear connection $``$ gives us the way to express Pontryagin characteristic classes of the vector bundle $``$ via the regular vector field $`\zeta `$. Indeed, any linear connection provides an expression for Pontryagin classes. We are going to briefly recall the correspondent classical construction (see for details). Let $`R^{}(\xi ,\eta )=[_\xi ,_\eta ]_{[\xi ,\eta ]}`$ be the curvature of linear connection $``$. Then $`R^{}(\xi ,\eta )\nu `$ is $`C^{\mathrm{}}(M)`$-linear with respect to each of three arguments $`\xi ,\eta ,\nu `$. In particular, $`R^{}(,)\nu (z)^2(T_z^{}N)T_zE_z,zN.`$ In other words, $`R^{}(,)\text{Hom}(,^2(T^{}N))`$.
Consider the commutative exterior algebra
$$^{ev}N=C^{\mathrm{}}(N)^2(T^{}N)\mathrm{}^{2n}(T^{}N)$$
of the even order differential forms on $`N`$. Then $`R^{}`$ can be treated as an endomorphism of the module $`^{ev}N`$ over algebra $`^{ev}N`$, i. e. $`R^{}\text{End}_{^{ev}N}\left(^{ev}M\right)`$. Now consider characteristic polynomial $`det(tI+\frac{1}{2\pi }R^{})=t^n+\underset{i=1}{\overset{n}{}}\varphi _it^{ni}`$, where the coefficient $`\varphi _i`$ is an order $`2i`$ differential form on $`N`$. All forms $`\varphi _i`$ are closed; the forms $`\varphi _{2k1}`$ are exact and the forms $`\varphi _{2k}`$ represent the Pontryagin characteristic classes, $`k=1,\mathrm{},[\frac{n}{2}]`$.
### 17 Symplectic setting
Assume that $`N`$ is a symplectic manifold endowed with a symplectic form $`\sigma `$. Recall that a symplectic form is just a closed non degenerate differential 2-form. Suppose $`E`$ is a Lagrange foliation on the symplectic manifold $`(N,\sigma )`$; this means that $`\sigma |_{E_z}=0,zN`$. Basic examples are cotangent bundles endowed with the standard symplectic structure: $`N=T^{}M,E_z=T_{\pi (z)}^{}M`$, where $`\pi :T^{}MM`$ is the canonical projection. In this case $`\sigma =d\tau `$, where $`\tau =\{\tau _z:zT^{}M\}`$ is the Liouville 1-form on $`T^{}M`$ defined by the formula: $`\tau _z=z\pi _{}`$. Completely integrable Hamiltonian systems provide another important class of Lagrange foliations. Weโll briefly recall the correspondent terminology. Details can be found in any introduction to symplectic geometry (for instance, in ).
Smooth functions on the symplectic manifold are called Hamiltonians. To any Hamiltonian there corresponds a Hamiltonian vector field $`\stackrel{}{h}`$ on $`M`$ defined by the equation: $`dh=\sigma (,\stackrel{}{h})`$. The Poisson bracket $`\{h_1,h_2\}`$ of the Hamiltonians $`h_1`$ and $`h_2`$ is the Hamiltonian defined by the formula: $`\{h_1,h_2\}=\sigma (\stackrel{}{h}_1,\stackrel{}{h}_2)=\stackrel{}{h}_1h_2`$. Poisson bracket is obviously anti-symmetric and satisfies the Jacobi identity: $`\{h_1,\{h_2,h_3\}\}+\{h_3,\{h_1,h_2\}\}+\{h_2,\{h_3,h_1\}\}=0`$. This identity is another way to say that the form $`\sigma `$ is closed. Jacobi identity implies one more useful formula: $`\stackrel{}{\{h_1,h_2\}}=[\stackrel{}{h}_1,\stackrel{}{h}_2]`$.
We say that Hamiltonians $`h_1,\mathrm{},h_n`$ are in involution if $`\{h_i,h_j\}=0`$; then $`h_j`$ is constant along trajectories of the Hamiltonian equation $`\dot{z}=\stackrel{}{h}_i(z),i,j=1,\mathrm{},n`$. We say that $`h_1,\mathrm{},h_n`$ are independent if $`d_zh_1\mathrm{}d_zh_n0,zN`$. $`n`$ independent Hamiltonians in involution form a completely integrable system. More precisely, any of Hamiltonian equations $`\dot{z}=\stackrel{}{h}_i(z)`$ is completely integrable with first integrals $`h_1,\mathrm{},h_n`$.
###### Lemma II.6
Let Hamiltonians $`h_1,\mathrm{},h_n`$ form a completely integrable system. Then the $`n`$-foliation $`E_z=\{z^{}M:h_i(z^{})=h_i(z),i=1,\mathrm{},n\},zN`$, is Lagrangian.
Proof. We have $`\stackrel{}{h}_ih_j=0,i,j=1,\mathrm{},n`$, hence $`\stackrel{}{h}_i(z)`$ are tangent to $`E_z`$. Vectors $`\stackrel{}{h}_1(z),\mathrm{},\stackrel{}{h}_n(z)`$ are linearly independent, hence
$$span\{\stackrel{}{h}_1(z),\mathrm{},\stackrel{}{h}_n(z)\}=T_zE_z.$$
Moreover, $`\sigma (\stackrel{}{h}_i,\stackrel{}{h}_j)=\{h_i,h_j\}=0`$, hence $`\sigma |_{E_z}=0.\mathrm{}`$
Any Lagrange foliation possesses a canonical affine structure. Let $`\left[\right]`$ be the sheaf of germs of the distribution $`=\{T_zE_z:zN\}`$ as in Section 16; then $`\left[\right]^a`$ is the intersection of $`\left[\right]`$ with the sheaf of germs of Hamiltonian vector fields.
We have to check that Lie algebra $`\left[\right]_z^a`$ is Abelian and generates $`T_zE_z,zN`$. First check the Abelian property. Let $`\stackrel{}{h}_1,\stackrel{}{h}_2\left[\right]_z^a`$; we have $`[\stackrel{}{h}_1,\stackrel{}{h}_2]=\stackrel{}{\{h_1,h_2\}},\{h_1,h_2\}=\sigma (\stackrel{}{h}_1,\stackrel{}{h}_2)=0`$, since $`\stackrel{}{h}_i`$ are tangent to $`E_z`$ and $`\sigma |_{E_z}=0`$. The second property follows from the DarbouxโWeinstein theorem (see ) which states that all Lagrange foliations are locally equivalent. More precisely, this theorem states that any $`zM`$ possesses a neighborhood $`O_z`$ and local coordinates which turn the restriction of the Lagrange foliation $`E`$ to $`O_z`$ into the trivial bundle $`^n\times ^n=\{(x,y):x,y^n\}`$ and, simultaneously, turn $`\sigma |_{O_z}`$ into the form $`\underset{i=1}{\overset{n}{}}dx_idy_i`$. In this special coordinates, the fibers become coordinate subspaces $`^n\times \{y\},y^n`$, and the required property is obvious: vector fields $`\frac{}{x_i}`$ are Hamiltonian fields associated to the Hamiltonians $`y_i,i=1,\mathrm{},n`$.
Suppose $`\zeta `$ is a Hamiltonian field on the symplectic manifold endowed with the Lagrange foliation, $`\zeta =\stackrel{}{h}`$. Let $`\varsigma \left[\right]_z^a,\varsigma =\stackrel{}{s}`$; then $`\varsigma h=\{s,h\}`$. The field $`\stackrel{}{h}`$ is regular if and only if the quadratic form $`s\{s,\{s,h\}\}(z)`$ has rank $`n`$. Indeed, in the โDarbouxโWeinstein coordinatesโ this quadratic form has the matrix $`\{\frac{^2h}{x_ix_j}\}_{i,j=1}^n`$.
Recall that the tangent space $`T_zN`$ to the symplectic manifold $`N`$ is a symplectic space endowed with the symplectic structure $`\sigma _z`$. An $`n`$-dimensional subspace $`\upsilon T_zN`$ is a Lagrangian subspace if $`\sigma _z|_\upsilon =0`$. The set
$$L(T_zN)=\{\upsilon G_n(T_zM):\sigma _z|_\upsilon =0\}$$
of all Lagrange subspaces of $`T_zM`$ is a Lagrange Grassmannian.
Hamiltonian flow $`e^{t\stackrel{}{h}}`$ preserves the symplectic form, $`\left(e^{t\stackrel{}{h}}\right)^{}\sigma =\sigma `$. Hence $`\left(e^{t\stackrel{}{h}}\right)_{}:T_zNT_{e^{t\stackrel{}{h}}(z)}N`$ transforms Lagrangian subspaces in the Lagrangian ones. It follows that the Jacobi curve $`J_z(t)=\left(e^{t\stackrel{}{h}}\right)_{}T_{e^{t\stackrel{}{h}}(z)}E_{e^{t\stackrel{}{h}}(z)}`$ consists of Lagrangian subspaces, $`J_z(t)L(T_zN)`$.
We need few simple facts on Lagrangian Grassmannians (see Sec. 6 for the basic information and \[3, Sec. 4\] for a consistent description of their geometry). Let $`(\mathrm{\Sigma },\overline{\sigma })`$ be a $`2n`$-dimensional symplectic space and $`\upsilon _0,\upsilon _1L(\mathrm{\Sigma })`$ be a pair of transversal Lagrangian subspaces, $`\upsilon _0\upsilon _1=0`$. Bilinear form $`\overline{\sigma }`$ induces a non degenerate pairing of the spaces $`\upsilon _0`$ and $`\upsilon _1`$ by the rule $`(e,f)\overline{\sigma }(e,f),e\upsilon _0,f\upsilon _1`$. To any basis $`e_1,\mathrm{},e_n`$ of $`\upsilon _0`$ we may associate a unique dual basis $`f_1,\mathrm{},f_n`$ of $`\upsilon _1`$ such that $`\overline{\sigma }(e_i,f_j)=\delta _{ij}`$. The form $`\overline{\sigma }`$ is totally normalized in the basis $`e_1,\mathrm{},e_n,f_1,\mathrm{},f_n`$ of $`\mathrm{\Sigma }`$, since $`\sigma (e_i,e_j)=\sigma (f_i,f_j)=0`$. It follows that symplectic group
$$\text{Sp}(\mathrm{\Sigma })=\{A\text{GL}(\mathrm{\Sigma }):\overline{\sigma }(Ae,Af)=\overline{\sigma }(e,f),e,f\mathrm{\Sigma }\}$$
acts transitively on the pairs of transversal Lagrangian subspaces.
Next result is a โsymplectic specificationโ of Lemma II.1 from Section 9.
###### Lemma II.7
Let $`\upsilon _0L(\mathrm{\Sigma })`$; then $`\{\pi _{\upsilon \upsilon _0}:\upsilon \upsilon _0^{}L(\mathrm{\Sigma })\}`$ is an affine subspace of the affine space $`\{\pi _{v\upsilon _0}:v\upsilon _0^{}\}`$ characterized by the relation:
$$v\upsilon _0^{}L(\mathrm{\Sigma })\overline{\sigma }(\pi _{v\upsilon _0},)+\overline{\sigma }(,\pi _{v\upsilon _0})=\overline{\sigma }(,).$$
Proof. Assume that $`\upsilon _1\upsilon _0^{}L(\mathrm{\Sigma })`$. Let $`e,f\mathrm{\Sigma },e=e_0+e_1,f=f_0+f_1`$ where $`e_i,f_i\upsilon _i,i=0,1`$; then
$$\overline{\sigma }(e,f)=\overline{\sigma }(e_0+e_1,f_0+f_1)=\overline{\sigma }(e_0,f_1)+\overline{\sigma }(e_1,f_0)=$$
$$\overline{\sigma }(e_0,f)+\overline{\sigma }(e,f_0)=\overline{\sigma }(\pi _{\upsilon _1\upsilon _0}e,f)+\overline{\sigma }(e,\pi _{\upsilon _1\upsilon _0}f).$$
Conversely, let $`v\upsilon _0^{}`$ is not a Lagrangian subspace. Then there exist $`e,fv`$ such that $`\overline{\sigma }(e,f)0`$, while $`\overline{\sigma }(\pi _{v\upsilon _0}e,f)=\overline{\sigma }(e,\pi _{v\upsilon _0}f)=0.\mathrm{}`$
###### Corollary II.4
Let $`v()`$ be an ample curve in $`G_n(\mathrm{\Sigma })`$ and $`v^{}()`$ be the derivative curve of $`v()`$. If $`v(t)L(\mathrm{\Sigma }),t`$, then $`v^{}(t)L(\mathrm{\Sigma })`$.
Proof. The derivative curve $`v^{}`$ was defined in Section 11. Recall that $`\pi _{v^{}(t)v(t)}=\pi _t^0`$, where $`\pi _t^0`$ is the free term of the Laurent expansion
$$\pi _{v(\tau )v(t)}\underset{i=k_t}{\overset{\mathrm{}}{}}(\tau t)^i\pi _t^i.$$
The free term $`\pi _t^0`$ belongs to the affine hull of $`\pi _{v(\tau )v(t)}`$, when $`\tau `$ runs a neighborhood of $`t`$. Since $`\pi _{v(\tau )v(t)}`$ belongs to the affine space $`\{\pi _{vv_0}:vv_0^{}L(\mathrm{\Sigma })\}`$, then $`\pi _t^0`$ belongs to this affine space as well. $`\mathrm{}`$
We call a Lagrange distribution any rank $`n`$ vector distribution $`\{\mathrm{\Lambda }_zT_zN:zN\}`$ on the symplectic manifold $`N`$ such that $`\mathrm{\Lambda }_zL(T_zN),zN`$.
###### Corollary II.5
Canonical Ehresmann connection $`_\zeta =\{J_z^{}(0):zN\}`$ associated to an ample Hamiltonian field $`\zeta =\stackrel{}{h}`$ is a Lagrange distribution. $`\mathrm{}`$
It is clearly seeing in coordinates how Lagrange Grassmanian is sitting in the usual one. Let $`\mathrm{\Sigma }=^n\times ^n=\{(\eta ,y):\eta ^n,y^n\}`$. Then any $`v\left(\{0\}\times ^n\right)^{}`$ has a form $`v=\{(y^{},Sy):y^n\}`$, where $`S`$ is an $`n\times n`$-matrix. It is easy to see that $`v`$ is a Lagrangian subspace if and only if $`S`$ is a symmetric matrix, $`S=S^{}`$.
### 18 Monotonicity
We continue to study curves in the Lagrange Grassmannian $`L(T_zN)`$, in particular, the Jacobi curves $`t\left(e^{t\stackrel{}{H}}\right)_{}T_{e^{t\stackrel{}{H}}(z)}E_{e^{t\stackrel{}{H}}(z)}`$. In Section 6 we identified the velocity $`\dot{\mathrm{\Lambda }}(t)`$ of any smooth curve $`\mathrm{\Lambda }()`$ in $`L(T_zN)`$ with a quadratic form $`\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}(t)`$ on the subspace $`\mathrm{\Lambda }(t)T_zN`$. Recall that the curve $`\mathrm{\Lambda }()`$ was called monotone increasing if $`\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}(t)0,t`$; it is called monotone decreasing if $`\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}(t)0`$. It is called monotone in both cases.
###### Proposition II.8
Set $`\mathrm{\Lambda }(t)=\left(e^{t\stackrel{}{H}}\right)_{}T_{e^{t\stackrel{}{H}}(z)}E_{e^{t\stackrel{}{H}}(z)}`$; then quadratic form $`\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}(t)`$ is equivalent (up to a linear change of variables) to the form
$$\varsigma (\varsigma \varsigma H)(e^{t\stackrel{}{H}}(z)),\varsigma []_{e^{t\stackrel{}{H}}(z)}^a,$$
$`(21)`$
on $`E_{e^{t\stackrel{}{H}}(z)}`$.
Proof. Let $`z_t=e^{t\stackrel{}{H}}(z)`$, then
$$\frac{d}{dt}\mathrm{\Lambda }(t)=\frac{d}{dt}e_{}^{(t_0t)\stackrel{}{H}}T_{z_t}E_{z_t}=e_{}^{(t_0t)\stackrel{}{H}}\frac{d}{d\epsilon }|_{\epsilon =0}e_{}^{\epsilon \stackrel{}{H}}T_{z_{t+\epsilon }}E_{z_{t+\epsilon }}.$$
Set $`\mathrm{\Delta }(\epsilon )=e_{}^{\epsilon \stackrel{}{H}}T_{z_{t+\epsilon }}E_{z_{t+\epsilon }}L\left(T_{z_t}N\right)`$. It is enough to prove that $`\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Delta }}}(0)`$ is equivalent to form (21). Indeed, $`\dot{\mathrm{\Lambda }}(t)=e_{}^{(t_0t)\stackrel{}{H}}T_{z_t}\dot{\mathrm{\Delta }}(0)`$, where
$$e_{}^{(t_0t)\stackrel{}{H}}:T_{z_t}NT_{z_{t_0}}N$$
is a symplectic isomorphism. The association of the quadratic form $`\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}(t)`$ on the subspace $`\mathrm{\Lambda }(t)`$ to the tangent vector $`\dot{\mathrm{\Lambda }}(t)L\left(T_{z_{t_0}}N\right)`$ is intrinsic, i.e. depends only on the symplectic structure on $`T_{z_{t_0}}N`$. Hence $`\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Delta }}}(0)(\xi )=\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}(t)\left(e_{}^{(t_0t)\stackrel{}{H}}\xi \right)`$, $`\xi \mathrm{\Delta }(0)=T_{z_t}E_{z_t}`$.
What remains, is to compute $`\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Delta }}}(0)`$; we do it in the DarbouxโWeinstein coordinates $`z=(x,y)`$. We have: $`\mathrm{\Delta }(\epsilon )=`$
$$\{(\xi (\epsilon ),\eta (\epsilon )):\begin{array}{ccc}\hfill \dot{\xi }(\tau )& =& \xi (\tau )\frac{^2H}{xy}(z_{t\tau })+\eta (\tau )^{}\frac{^2H}{y^2}(z_{t\tau }),\hfill \\ \hfill \dot{\eta }(\tau )& =& \frac{^2H}{x^2}(z_{t\tau })\xi (\tau )^{}\frac{^2H}{yx}(z_{t\tau })\eta (\tau ),\hfill \end{array}\genfrac{}{}{0pt}{}{\xi (0)=\xi ^n}{\eta (0)=0^n}\},$$
$$\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Delta }}}(0)(\xi )=\sigma ((\xi ,0),(\dot{\xi }(0),\dot{\eta }(0)))=\xi \dot{\eta }(0)=\xi \frac{^2H}{x^2}(z_t)\xi ^{}.$$
Recall now that form (21) has matrix $`\frac{^2H}{x^2}(z_t)`$ in the DarbouxโWeinstein coordinates. $`\mathrm{}`$
This proposition clearly demonstrates the importance of monotone curves. Indeed, monotonicity of Jacobi curves is equivalent to the convexity (or concavity) of the Hamiltonian on each leaf of the Lagrange foliation. In the case of a cotangent bundle this means the convexity or concavity of the Hamiltonian with respect to the impulses. All Hamiltonians (energy functions) of mechanical systems are like that! This is not an occasional fact but a corollary of the list action principle. Indeed, trajectories of the mechanical Hamiltonian system are extremals of the least action principle and the energy function itself is the Hamiltonian of the correspondent regular optimal control problem as it was considered in Section 7. Moreover, it was stated in Section 7 that convexity of the Hamiltonian with respect to the impulses is necessary for the extremals to have finite Morse index. It turns out that the relation between finiteness of the Morse index and monotonicity of the Jacobi curve has a fundamental nature. A similar property is valid for any, not necessary regular, extremal of a finite Morse index. Of course, to formulate this property we have first to explain what are Jacobi curve for non regular extremals. To do that, we come back to the very beginning; indeed, Jacobi curves appeared first as the result of calculation of the $``$-derivative at the regular extremal (see Sections 7, 8). On the other hand, $``$-derivative is well-defined for any extremal of the finite Morse index as it follows from Theorem I.1. We thus come to the following construction in which we use notations and definitions of Sections 3, 4.
Let $`h(\lambda ,u)`$ be the Hamiltonian of a smooth optimal control system, $`\lambda _t,t_0tt_1`$, an extremal, and $`q(t)=\pi (\lambda _t),t_0,tt_1`$ the extremal path. Recall that the pair $`(\lambda _{t_0},\lambda _t)`$ is a Lagrangian multiplier for the conditional minimum problem defined on an open subset of the space
$$M\times L_{\mathrm{}}([t_0,t_1],U)=\{(q_t,u()):qM,u()L_{\mathrm{}}([t_0,t_1],U)\},$$
where $`u()`$ is control and $`q_t`$ is the value at $`t`$ of the solution to the differential equation $`\dot{q}=f(q,u(\tau )),\tau [t_0,t_1]`$. In particular, $`F_t(q_t,u())=q_t`$. The cost is $`J_{t_0}^{t_1}(q_t,u())`$ and constraints are $`F_{t_0}(q_t,u())=q(0),q_t=q(t)`$.
Let us set $`J_t(u)=J_{t_0}^t(q(t),u()),\mathrm{\Phi }_t(u)=F_{t_0}(q(t),u())`$. A covector $`\lambda T^{}M`$ is a Lagrange multiplier for the problem $`(J_t,\mathrm{\Phi }_t)`$ if and only if there exists an extremal $`\widehat{\lambda }_\tau ,t_0\tau t`$, such that $`\lambda _{t_0}=\lambda ,\widehat{\lambda }_tT_{q(t)}^{}M`$. In particular, $`\lambda _{t_0}`$ is a Lagrange multiplier for the problem $`(J_t,\mathrm{\Phi }_t)`$ associated to the control $`u()=\overline{u}(\lambda _.)`$.
Assume that $`\mathrm{ind}\mathrm{Hess}_u\left(J_{t_1}|_{\mathrm{\Phi }_{t_1}^1(q(t_0))}\right)\mathrm{},t_0tt_1`$ and set $`\overline{\mathrm{\Phi }}_t=(J_t,\mathrm{\Phi }_t)`$. The curve
$$t_{(\lambda _{t_0},u)}(\overline{\mathrm{\Phi }}_t),t_0tt_1$$
in the Lagrange Grassmannian $`L\left(T_{\lambda _{t_0}}(T^{}M)\right)`$ is called the Jacobi curve associated to the extremal $`\lambda _t,t_0tt_1`$.
In general, the Jacobi curve $`t_{(\lambda _{t_0},u)}(\overline{\mathrm{\Phi }}_t)`$ is not smooth, it may even be discontinues, but it is monotone decreasing in a sense we are going to briefly describe now. You can find more details in (just keep in mind that similar quantities may have opposite signs in different papers; sign agreements vary from paper to paper that is usual for symplectic geometry). Monotone curves in the Lagrange Grassmannian have analytic properties similar to scalar monotone functions: no more than a countable set of discontinuity points, right and left limits at every point, and differentiability almost everywhere with semi-definite derivatives (nonnegative for monotone increasing curves and nonpositive for decreasing ones). True reason for such a monotonicity is a natural monotonicity of the family $`\overline{\mathrm{\Phi }}_t`$. Indeed, let $`\tau <t`$, then $`\overline{\mathrm{\Phi }}_\tau `$ is, in fact, the restriction of $`\overline{\mathrm{\Phi }}_t`$ to certain subspace: $`\overline{\mathrm{\Phi }}_\tau =\overline{\mathrm{\Phi }}_t๐ญ_\tau `$, where $`๐ญ_\tau (u)(s)=\{\begin{array}{ccc}\hfill u(s)& ,& s<\tau \hfill \\ \hfill \stackrel{~}{u}(s)& ,& s>\tau \hfill \end{array}`$. One can define the Maslov index of a (maybe discontinues) monotone curve in the Lagrange Grassmannian and the relation between the Morse and Maslov index indices from Theorem I.3 remains true.
In fact, Maslov index is a key tool in the whole construction. The starting point is the notion of a simple curve. A smooth curve $`\mathrm{\Lambda }(\tau ),\tau _0\tau \tau _1,`$ in the Lagrange Grassmannian $`L(\mathrm{\Sigma })`$ is called simple if there exists $`\mathrm{\Delta }L(\mathrm{\Sigma })`$ such that $`\mathrm{\Delta }\mathrm{\Lambda }(\tau )=0,\tau [\tau _0,\tau _1]`$; in other words, the entire curve is contained in one coordinate chart. It is not hard to show that any two points of $`L(\mathrm{\Sigma })`$ can be connected by a simple monotone increasing (as well as monotone decreasing) curve. An important fact is that the Maslov index $`\mu (\mathrm{\Lambda }_\mathrm{\Pi }())`$ of a simple monotone increasing curve $`\mathrm{\Lambda }(\tau ),\tau _0\tau \tau _1`$ is uniquely determined by the triple $`(\mathrm{\Pi },\mathrm{\Lambda }(\tau _0),\mathrm{\Lambda }(\tau _1))`$; i.e. it has the same value for all simple monotone increasing curves connecting $`\mathrm{\Lambda }(\tau _0)`$ with $`\mathrm{\Lambda }(\tau _1)`$. A simple way to see this is to find an intrinsic algebraic expression for the Maslov index preliminary computed for some simple monotone curve in some coordinates. We can use Lemma I.2 for this computation since the curve is simple. The monotonic increase of the curve implies that $`S_{\mathrm{\Lambda }(t_1)}>S_{\mathrm{\Lambda }(t_0)}`$.
Exercise. Let $`S_0,S_1`$ be nondegenerate symmetric matrices and $`S_1S_0`$. Then $`\mathrm{ind}S_0\mathrm{ind}S_1=\mathrm{ind}\left(S_0^1S_1^1\right).`$
Let $`x(\mathrm{\Lambda }(\tau _0)+\mathrm{\Lambda }(\tau _1)\mathrm{\Pi }`$ so that $`x=x_0+x_1`$, where $`x_i\mathrm{\Lambda }(\tau _i),i=0,1`$. We set $`๐ฎ(x)=\sigma (x_1,x_0)`$. If $`\mathrm{\Lambda }(\tau _0)\mathrm{\Lambda }(\tau _1)=0`$, then $`\mathrm{\Lambda }(\tau _0)+\mathrm{\Lambda }(\tau _1)=\mathrm{\Sigma }`$, $`x`$ is any element of $`\mathrm{\Pi }`$ and $`x_0,x_1`$ are uniquely determined by $`x`$. This is not true if $`\mathrm{\Lambda }(\tau _0)\mathrm{\Lambda }(\tau _1)0`$ but $`๐ฎ(x)`$ is well-defined anyway: $`\sigma (x_1,x_2)`$ depends only on $`x_0+x_1`$ since $`\sigma `$ vanishes on $`\mathrm{\Lambda }(\tau _i),i=0,1.`$
Now we compute $`๐ฎ`$ in coordinates. Recall that
$$\mathrm{\Lambda }(\tau _i)=\{(y^{},S_{\mathrm{\Lambda }(\tau _i)}y):y^n\},i=0,1,\mathrm{\Pi }=\{y^{},0):y^n\}.$$
We have
$$๐ฎ(x)=y_1^{}S_{\mathrm{\Lambda }(\tau _0)}y_0y_0^{}S_{\mathrm{\Lambda }(\tau _1)}y_1,$$
where $`x=(y_0^{}+y_1^{},0)`$, $`S_{\mathrm{\Lambda }(\tau _0)}y_0+S_{\mathrm{\Lambda }(\tau _1)}y_1=0`$. Hence $`y_1=S_{\mathrm{\Lambda }(tau_1)}^1S_{\mathrm{\Lambda }(\tau _0)}y_0`$ and
$$๐ฎ(x)=y_0^{}S_{\mathrm{\Lambda }(\tau _0)}y_0\left(S_{\mathrm{\Lambda }(\tau _0)}y_0\right)^{}S_{\mathrm{\Lambda }(\tau _1)}^1S_{\mathrm{\Lambda }(\tau _0)}y_0=y^{}\left(S_{\mathrm{\Lambda }(\tau _0)}^1S_{\mathrm{\Lambda }(\tau _1)}^1\right)y,$$
where $`y=S_{\mathrm{\Lambda }(\tau _0)}y_0`$. We see that the form $`๐ฎ`$ is equivalent, up to a linear change of coordinates, to the quadratic form defined by the matrix $`S_{\mathrm{\Lambda }(\tau _0)}^1S_{\mathrm{\Lambda }(\tau _1)}^1`$. Now we set
$$\mathrm{ind}_\mathrm{\Pi }(\mathrm{\Lambda }(\tau _0),\mathrm{\Lambda }(\tau _1))\stackrel{def}{=}\mathrm{ind}๐ฎ.$$
The above exercise and Lemma I.2 imply the following:
###### Lemma II.8
If $`\mathrm{\Lambda }(\tau ),\tau _0\tau \tau _1`$, is a simple monotone increasing curve, then
$$\mu (\mathrm{\Lambda }())=\mathrm{ind}_\mathrm{\Pi }(\mathrm{\Lambda }(\tau _0),\mathrm{\Lambda }(\tau _1)).$$
Note that definition of the form $`๐ฎ`$ does not require transversality of $`\mathrm{\Lambda }(\tau _i)`$ to $`\mathrm{\Pi }`$. It is convenient to extend definition of $`\mathrm{ind}_\mathrm{\Pi }(\mathrm{\Lambda }(\tau _0),\mathrm{\Lambda }(\tau _1))`$ to this case. General definition is as follows:
$$\mathrm{ind}_\mathrm{\Pi }(\mathrm{\Lambda }_0,\mathrm{\Lambda }_1)=\mathrm{ind}๐ฎ+\frac{1}{2}(dim(\mathrm{\Pi }\mathrm{\Lambda }_0)+dim(\mathrm{\Pi }\mathrm{\Lambda }_1))dim(\mathrm{\Pi }\mathrm{\Lambda }_0\mathrm{\Lambda }_1).$$
The Maslov index also has appropriate extension (see \[3, Sec.4\]) and Lemma II.8 remains true.
Index $`\mathrm{ind}_\mathrm{\Pi }(\mathrm{\Lambda }_0,\mathrm{\Lambda }_1)`$ satisfies the triangle inequality:
$$\mathrm{ind}_\mathrm{\Pi }(\mathrm{\Lambda }_0,\mathrm{\Lambda }_2)\mathrm{ind}_\mathrm{\Pi }(\mathrm{\Lambda }_0,\mathrm{\Lambda }_1)+\mathrm{ind}_\mathrm{\Pi }(\mathrm{\Lambda }_1,\mathrm{\Lambda }_2).$$
Indeed, the right-hand side of the inequality is equal to the Maslov index of a monotone increasing curve connecting $`\mathrm{\Lambda }_0`$ with $`\mathrm{\Lambda }_2`$, i.e. of the concatenation of two simple monotone increasing curves. Obviously, the Maslov index of a simple monotone increasing curve is not greater than the Maslov index of any other monotone increasing curve connecting the same endpoints.
The constructed index gives a nice presentation of the Maslov index of any (not necessary simple) monotone increasing curve $`\mathrm{\Lambda }(t),t_0tt_1`$:
$$\mu _\mathrm{\Pi }(\mathrm{\Lambda }())=\underset{i=0}{\overset{l}{}}\mathrm{ind}_\mathrm{\Pi }(\mathrm{\Lambda }(\tau _i),\mathrm{\Lambda }(\tau _{i+1})),$$
$`(22)`$
where $`t_0=\tau _0<\tau _1<\mathrm{}<\tau _l<\tau _{l+1}=t_1`$ and $`\mathrm{\Lambda }|_{[\tau _i,\tau _{i+1}]}`$ are simple pieces of the curve $`\mathrm{\Lambda }()`$. If the pieces are not simple, then the right-hand side of (22) gives a low bound for the Maslov index (due to the triangle inequality).
Let now $`\mathrm{\Lambda }(t),t_0tt_1,`$ be a smooth curve which is not monotone increasing. Take any subdivision $`t_0=\tau _0<\tau _1<\mathrm{}<\tau _l<\tau _{l+1}=t_1`$ and compute the sum $`\underset{i=0}{\overset{l}{}}\mathrm{ind}_\mathrm{\Pi }(\mathrm{\Lambda }(\tau _i),\mathrm{\Lambda }(\tau _{i+1}))`$. This sum inevitably goes to infinity when the subdivision becomes finer and finer. The reason is as follows: $`\mathrm{ind}_\mathrm{\Pi }(\mathrm{\Lambda }(\tau _i),\mathrm{\Lambda }(\tau _{i+1}))>0`$ for any simple piece $`\mathrm{\Lambda }|_{[\tau _i,\tau _{i+1}]}`$ such that $`\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}(\tau )0,\tau [\tau _i,\tau _{i+1}]`$ and $`\mu _\mathrm{\Pi }(\mathrm{\Lambda }|_{[\tau _i,\tau _{i+1}]}=0`$. I advise reader to play with the one-dimensional case of the curve in $`L(^2)=S^1`$ to see better whatโs going on.
This should now be clear how to manage in the general nonsmooth case. Take a curve $`\mathrm{\Lambda }()`$ (an arbitrary mapping from $`[t_0,t_1]`$ into $`L(\mathrm{\Sigma })`$). For any finite subset $`๐ฏ+\{\tau _1,\mathrm{},\tau _k\}[t_0,t_1]`$, where $`t_0=\tau _0<\tau _1<\mathrm{}<\tau _l<\tau _{l+1}=t_1`$, we compute the sum $`I_\mathrm{\Pi }^๐ฏ=\underset{i=0}{\overset{l}{}}\mathrm{ind}_\mathrm{\Pi }(\mathrm{\Lambda }(\tau _i),\mathrm{\Lambda }(\tau _{i+1}))`$ and then find supremum of these sums for all finite subsets: $`I_\mathrm{\Pi }(\mathrm{\Lambda }())=\underset{๐ฏ}{sup}I_\mathrm{\Pi }^๐ฏ`$. The curve $`\mathrm{\Lambda }()`$ is called monotone increasing if $`I_\mathrm{\Pi }^๐ฏ<\mathrm{}`$; it is not hard to show that the last property does not depend on $`\mathrm{\Pi }`$ and that monotone increased curves enjoy listed above analytic properties. A curve $`\mathrm{\Lambda }()`$ is called monotone decreasing if inversion of the parameter $`tt_0+t_1t`$ makes it monotone increasing.
We set $`\mu (\mathrm{\Lambda }())=I_\mathrm{\Pi }(\mathrm{\Lambda }())`$ for any monotone increasing curve and $`\mu (\mathrm{\Lambda }())=I_\mathrm{\Pi }(\widehat{\mathrm{\Lambda }}())`$ for a monotone decreasing one, where $`\widehat{\mathrm{\Lambda }}(t)=\mathrm{\Lambda }(t_0+t_1t)`$. The defined in this way Maslov index of a discontinues monotone curve equals the Maslov index of the continues curve obtained by gluing all discontinuities with simple monotone curves of the same direction of monotonicity.
If $`\mathrm{\Lambda }(t)=_{(\lambda _{t_0},u)}(\overline{\mathrm{\Phi }}_t)`$ is the Jacobi curve associated to the extremal with a finite Morse index, then $`\mathrm{\Lambda }()`$ is monotone decreasing and its Maslov index computes $`\mathrm{ind}\mathrm{Hess}_u\left(J_{t_1}|_{\mathrm{\Phi }_{t_1}^1(q(t_0))}\right)`$ in the way similar to Theorem I.3. Of course, these nice things have some value only if we can effectively find Jacobi curves for singular extremals: their definition was too abstract. Fortunately, this is not so hard; see for the explicit expression of Jacobi curves for a wide class of singular extremals and, in particular, for singular curves of rank 2 vector distributions (these last Jacobi curves have found important applications in the geometry of distributions, see ).
One more important property of monotonic curves is as follows.
###### Lemma II.9
Assume that $`\mathrm{\Lambda }()`$ is monotone and right-continues at $`t_0`$, i.e. $`\mathrm{\Lambda }(t_0)=\underset{tt_0}{lim}\mathrm{\Lambda }(t)`$. Then $`\mathrm{\Lambda }(t_0)\mathrm{\Lambda }(t)=\underset{t_0\tau t}{}\mathrm{\Lambda }(t)`$ for any $`t`$ sufficiently close to (and greater than) $`t_0`$.
Proof. We may assume that $`\mathrm{\Lambda }()`$ is monotone increasing. Take centered at $`\mathrm{\Lambda }(t_0)`$ local coordinates in the Lagrange Grassmannian; the coordinate presentation of $`\mathrm{\Lambda }(t)`$ is a symmetric matrix $`S_{\mathrm{\Lambda }(t)}`$, where $`S_{\mathrm{\Lambda }(t_0)}=0`$ and $`ty^{}S_{\mathrm{\Lambda }(t)}y`$ is a monotone increasing scalar function $`y^n`$. In particular, $`\mathrm{ker}S_{\mathrm{\Lambda }(t)}=\mathrm{\Lambda }(t)\mathrm{\Lambda }(t_0)`$ is a monotone decreasing family of subspaces. $`\mathrm{}`$
We set $`\mathrm{\Gamma }_t=\underset{t_0\tau t}{}\mathrm{\Lambda }(\tau )`$, a monotone decreasing family of isotropic subspaces. Let $`\mathrm{\Gamma }=\underset{t>t_0}{\mathrm{max}}\mathrm{\Gamma }_t`$, then $`\mathrm{\Gamma }_t==\mathrm{\Gamma }`$ for all $`t>t_0`$ sufficiently close to $`t_0`$. We have: $`\mathrm{\Lambda }(t)=\mathrm{\Lambda }(t)^{\mathrm{}}`$ and $`\mathrm{\Lambda }(t)\mathrm{\Gamma }`$ for all $`t>t_0`$ close enough to $`t_0`$; hence $`\mathrm{\Gamma }_t^{\mathrm{}}\mathrm{\Lambda }(t)`$. In particular, $`\mathrm{\Lambda }(t)`$ can be treated as a Lagragian subspace of the symplectic space $`\mathrm{\Gamma }^{\mathrm{}}/\mathrm{\Gamma }`$. Moreover, Lemma II.9 implies that $`\mathrm{\Lambda }(t)\mathrm{\Lambda }(t_0)=\mathrm{\Gamma }`$. In other words, $`\mathrm{\Lambda }(t)`$ is transversal to $`\mathrm{\Lambda }(t_0)`$ in $`\mathrm{\Gamma }^{\mathrm{}}/\mathrm{\Gamma }`$. In the case of a real-analytic monotone curve $`\mathrm{\Lambda }()`$ this automatically implies that $`\mathrm{\Lambda }()`$ is an ample curve in $`\mathrm{\Gamma }^{\mathrm{}}/\mathrm{\Gamma }`$. Hence any nonconstant monotone analytic curve is reduced to an ample monotone curve. It becoms ample after the factorization by a fixed (motionless) subspace.
### 19 Comparizon theorem
We come back to smooth regular curves after the deviation devoted to a more general perspective.
###### Lemma II.10
Let $`\mathrm{\Lambda }(t),t[t_0,t_1]`$ be a regular monotone increasing curve in the Lagrange Grassmannian $`L(\mathrm{\Sigma })`$. Then $`\{t[t_0,t_1]:\mathrm{\Lambda }(t)\mathrm{\Pi }0\}`$ is a finite subset of $`[t_0,t_1]\mathrm{\Pi }L(\mathrm{\Sigma })`$. If $`t_0`$ and $`t_1`$ are out of this subset, then
$$\mu _\mathrm{\Pi }(\mathrm{\Lambda }())=\underset{t(t_0,t_1)}{}dim(\mathrm{\Lambda }(t)\mathrm{\Pi }).$$
Proof. We have to proof that $`\mathrm{\Lambda }(t)`$ may have a nontrivial intersection with $`\mathrm{\Pi }`$ only for isolated values of $`t`$; the rest is Lemma I.1. Assume that $`\mathrm{\Lambda }(t)\mathrm{\Pi }0`$ and take a centered at $`\mathrm{\Pi }`$ coordinate neighborhood in $`L(\mathrm{\Sigma })`$ which contains $`\mathrm{\Lambda }(t)`$. In these coordinates, $`\mathrm{\Lambda }(\tau )`$ is presented by a symmetric matrix $`S_\mathrm{\Lambda }(\tau )`$ for any $`\tau `$ sufficiently close to $`t`$ and $`\mathrm{\Lambda }(\tau )\mathrm{\Pi }=\mathrm{ker}S_{\mathrm{\Lambda }(\tau )}.`$ Monotonicity and regularity properties are equivalent to the inequality $`\dot{S}_{\mathrm{\Lambda }(\tau )}>0`$. In particular, $`y^{}\dot{S}_{\mathrm{\Lambda }(t)}y>0y\mathrm{ker}S_{\mathrm{\Lambda }(t)}\{0\}`$. The last inequality implies that $`S_{\mathrm{\Lambda }(\tau )}`$ is a nondegenerate for all $`\tau `$ sufficiently close and not equal to $`t`$.
Definition. Parameter values $`\tau _0,\tau _1`$ are called conjugate for the continues curve $`\mathrm{\Lambda }()`$ in the Lagrange Grassmannian if $`\mathrm{\Lambda }(\tau _0)\mathrm{\Lambda }(\tau _1)0`$; the dimension of $`\mathrm{\Lambda }(\tau _0)\mathrm{\Lambda }(\tau _1)`$ is the multiplicity of the conjugate parameters.
If $`\mathrm{\Lambda }()`$ is a regular monotone increasing curve, then, according to Lemma II.9, conjugate points are isolated and the Maslov index $`\mu _{\mathrm{\Lambda }(t_0)}\left(\mathrm{\Lambda }|_{[t,t_1]}\right)`$ equals the sum of multiplicities of the conjugate to $`t_0`$ parameter values located in $`(t,t_1)`$. If $`\mathrm{\Lambda }()`$ is the Jacobi curve of an extremal of an optimal control problem, then this Maslov index equals the Morse index of the extremal; this is why conjugate points are so important.
Given a regular monotone curve $`\mathrm{\Lambda }()`$, the quadratic form $`\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}(t)`$ defines an Euclidean structure $`,_{\dot{\mathrm{\Lambda }}(t)}`$ on $`\mathrm{\Lambda }(t)`$ so that $`\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}(t)(x)=x,x_{\dot{\mathrm{\Lambda }}(t)}`$. Let $`R_\mathrm{\Lambda }(t)\mathrm{gl}(\mathrm{\Lambda }(t))`$ be the curvature operator of the curve $`\mathrm{\Lambda }()`$; we define the curvature quadratic form $`r_\lambda (t)`$ on $`\mathrm{\Lambda }(t)`$ by the formula:
$$r_\mathrm{\Lambda }(t)(x)=R_\mathrm{\Lambda }(t)x,x_{\dot{\mathrm{\Lambda }}(t)},x\mathrm{\Lambda }(t).$$
###### Proposition II.9
The curvature operator $`R_\mathrm{\Lambda }(t)`$ is a self-adjoint operator for the Euclidean structure $`,_{\dot{\mathrm{\Lambda }}(t)}`$. The form $`r_\mathrm{\Lambda }(t)`$ is equivalent (up to linear changes of variables) to the form $`\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}^{}(t)`$, where $`\mathrm{\Lambda }^{}()`$ is the derivative curve.
Proof. The statement is intrinsic and we may check it in any coordinates. Fix $`t`$ and take Darboux coordinates $`\{(\eta ,y):\eta ^n,y^n\}`$ in $`\mathrm{\Sigma }`$ in such a way that $`\mathrm{\Lambda }(t)=\{(y^{},0):y^n\}`$, $`\mathrm{\Lambda }^{}(t)=\{(0,y):y^n\}`$, $`\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}(t)(y)=y^{}y`$. Let $`\mathrm{\Lambda }(\tau )=\{(y^{},S_\tau y):y^n\}`$, then $`S_t=0`$. Moreover, $`\dot{S}(t)`$ is the matrix of the form $`\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}(t)`$ in given coordinates, hence $`\dot{S}_t=I`$. Recall that $`\mathrm{\Lambda }^{}(\tau )=\{(y^{}A_\tau ,y+S_\tau A_\tau y):y^n\}`$, where $`A_\tau =\frac{1}{2}\dot{S}_\tau ^1\ddot{S}_\tau \dot{S}_\tau ^1`$ (see (5)). Hence $`\ddot{S}_t=0`$. We have: $`R_\mathrm{\Lambda }(t)=\frac{1}{2}\underset{t}{\overset{\mathrm{}}{S}}`$, $`r_\mathrm{\Lambda }(t)(y)=\frac{1}{2}y^{}\underset{t}{\overset{\mathrm{}}{S}}y`$,
$$\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}^{}(t)(y)=\sigma ((0,y),(y^{}\dot{A}_t,0))=y^{}\dot{A}_ty=\frac{1}{2}y^{}\underset{t}{\overset{\mathrm{}}{S}}y.$$
So $`r_\mathrm{\Lambda }(t)`$ and $`\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}^{}(t)`$ have equal matrices for our choice of coordinates in $`\mathrm{\Lambda }(t)`$ and $`\mathrm{\Lambda }^{}(t)`$. The curvature operator is self-adjoint since it is presented by a symmetric matrix in coordinates where form $`\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}(t)`$ is the standard inner product. $`\mathrm{}`$
Proposition II.9 implies that the curvature operators of regular monotone curves in the Lagrange Grassmannian are diagonalizable and have only real eigenvalues.
###### Theorem II.1
Let $`\mathrm{\Lambda }()`$ be a regular monotone curve in the Lagrange Grassmannian $`L(\mathrm{\Sigma })`$, where $`dim\mathrm{\Sigma }=2n`$.
* If all eigenvalues of $`R_\mathrm{\Lambda }(t)`$ do not exceed a constant $`c0`$ for any $`t`$ from the domain of $`\mathrm{\Lambda }()`$, then $`|\tau _1\tau _0|\frac{\pi }{\sqrt{c}}`$ for any pair of conjugate parameter values $`\tau _0,\tau _1`$. In particular, If all eigenvalues of $`R_\mathrm{\Lambda }(t)`$ are nonpositive $`t`$, then $`\mathrm{\Lambda }()`$ does not possess conjugate parameter values.
* If $`\mathrm{tr}R_\mathrm{\Lambda }(t)nc`$ for some constant $`c>0`$ and $`t`$, then, for arbitrary $`\tau _0t`$, the segment $`[t,t+\frac{\pi }{\sqrt{c}}]`$ contains a conjugate to $`\tau _0`$ parameter value as soon as this segment is contained in the domain of $`\mathrm{\Lambda }()`$.
Both estimates are sharp.
Proof. We may assume without lack of generality that $`\mathrm{\Lambda }()`$ is ample monotone increasing. We start with the case of nonpositive eigenvalues of $`R_\mathrm{\Lambda }(t)`$. The absence of conjugate points follows from Proposition II.9 and the following
###### Lemma II.11
Assume that $`\mathrm{\Lambda }()`$ is an ample monotone increasing (decreasing) curve and $`\mathrm{\Lambda }^{}()`$ is a continues monotone decreasing (increasing) curve. Then $`\mathrm{\Lambda }()`$ does not possess conjugate parameter values and there exists a $`\underset{t+\mathrm{}}{lim}\mathrm{\Lambda }(t)=\mathrm{\Lambda }_{\mathrm{}}`$.
Proof. Take some value of the parameter $`\tau _0`$; then $`\mathrm{\Lambda }(\tau _0)`$ and $`\mathrm{\Lambda }^{}(\tau _0)`$ is a pair of transversal Lagrangian subspaces. We may choose coordinates in the Lagrange Grassmannian in such a way that $`S_{\mathrm{\Lambda }(\tau ))}=0`$ and $`S_{\mathrm{\Lambda }^{}(\tau _0)}=I`$, i.e. $`\mathrm{\Lambda }(\tau _0)`$ is represented by zero $`n\times n`$-matrix and $`\mathrm{\Lambda }^{}(\tau _0)`$ by the unit matrix. Monotonicity assumption implies that $`tS_{\mathrm{\Lambda }(t)}`$ is a monotone increasing curve in the space of symmetric matrices and $`tS_{\mathrm{\Lambda }^{}(t)}`$ is a monotone decreasing curve. Moreover, transversality of $`\mathrm{\Lambda }(t)`$ and $`\mathrm{\Lambda }^{}(t)`$ implies that $`S_{\mathrm{\Lambda }^{}(t)}S_{\mathrm{\Lambda }(t)}`$ is a nondegenerate matrix. Hence $`0<S_{\mathrm{\Lambda }(t)}<S_{\mathrm{\Lambda }^{}(t)}I`$ for any $`t>\tau _0`$. In particular, $`\mathrm{\Lambda }(t)`$ never leaves the coordinate neighborhood under consideration for $`T>\tau _0`$, the subspace $`\mathrm{\Lambda }(t)`$ is always transversal to $`\mathrm{\Lambda }(\tau _0)`$ and has a limit $`\mathrm{\Lambda }_{\mathrm{}}`$, where $`S_\mathrm{\Lambda }_{\mathrm{}}=\underset{t\tau _0}{sup}S_{\mathrm{\Lambda }(t)}.\mathrm{}`$
Now assume that the eigenvalues of $`R_\mathrm{\Lambda }(t)`$ do not exceed a constant $`c>0`$. We are going to reparametrize the the curve $`\mathrm{\Lambda }()`$ and to use the chain rule (7). Take some $`\overline{t}`$ in the domain of $`\mathrm{\Lambda }()`$ and set
$$\phi (t)=\frac{1}{\sqrt{c}}\left(\mathrm{arctan}(\sqrt{c}t)+\frac{\pi }{2}\right)+\overline{t},\mathrm{\Lambda }_\phi (t)=\mathrm{\Lambda }(\phi (t)).$$
We have: $`\phi ()=(\overline{t},\overline{t}+\frac{\pi }{\sqrt{c}})`$, $`\dot{\phi }(t)=\frac{1}{ct^2+1}`$, $`R_\phi (t)=\frac{c}{(ct^2+1)^2}`$. Hence, according to the chain rule (7), the operator
$$R_{\mathrm{\Lambda }_\phi }(t)=\frac{1}{(ct^2+1)^2}\left(R_\mathrm{\Lambda }(\phi (t))cI\right)$$
has only nonpositive eigenvalues. Already proved part of the theorem implies that $`\mathrm{\Lambda }_\phi `$ does not possess conjugate values of the parameter. In other words, any length $`\frac{\pi }{\sqrt{c}}`$ interval in the domain of $`\mathrm{\Lambda }()`$ is free of conjugate pairs of the parameter values.
Assume now that $`\mathrm{tr}R_\mathrm{\Lambda }(t)nc`$. We will prove that the existence of $`\mathrm{\Delta }L(\mathrm{\Sigma })`$ such that $`\mathrm{\Delta }\mathrm{\Lambda }(t)=0`$ for all $`t[\overline{t},\tau ]`$ implies that $`\tau \overline{t}<\frac{\pi }{\sqrt{c}}`$. Weโll prove it by contradiction. If there exists such a $`\mathrm{\Delta }`$, then $`\mathrm{\Lambda }|_{[\overline{t},\tau ]}`$ is completely contained in a fixed coordinate neighborhood of $`L(\mathrm{\Sigma })`$, therefore the curvature operator $`R_\mathrm{\Lambda }(t)`$ is defined by the formula (6). Put $`B(t)=(2\dot{S}_t)^1\ddot{S}_t`$, $`b(t)=\mathrm{tr}B(t)`$, $`t[\overline{t},\tau ]`$. Then
$$\dot{B}(t)=B^2(t)+R_\mathrm{\Lambda }(t),\dot{b}(t)=\mathrm{tr}B^2(t)+\mathrm{tr}R_\mathrm{\Lambda }(t).$$
Since for an arbitrary symmetric $`n\times n`$-matrix $`A`$ we have $`\mathrm{tr}A^2\frac{1}{n}(\mathrm{tr}A)^2`$, the inequality $`\dot{b}\frac{b^2}{n}+nc`$ holds. Hence $`b(t)\beta (t),\overline{t}t\tau `$, where $`\beta ()`$ is a solution of the equation $`\dot{\beta }=\frac{\beta ^2}{n}+nc`$, i.e. $`\beta (t)=n\sqrt{c}\mathrm{tan}(\sqrt{c}(tt_0))`$. The function $`b()`$ together with $`\beta ()`$ are bounded on the segment $`[\overline{t},\tau ]`$. Hence $`\tau t\frac{\pi }{\sqrt{c}}`$.
To verify that the estimates are sharp, it is enough to consider regular monotone curves of constant curvature. $`\mathrm{}`$
### 20 Reduction
We consider a Hamiltonian system on a symplectic manifold $`N`$ endowed with a fixed Lagrange foliation $`E`$. Assume that $`g:N`$ is a first integral of our Hamiltonian system, i.e. $`\{h,g\}=0`$.
###### Lemma II.12
Let $`zN,g(z)=c`$. The leaf $`E_z`$ is transversal to $`g^1(c)`$ at $`z`$ if and only if $`\stackrel{}{g}(z)T_zE_z`$.
Proof. Hypersurface $`g^1(c)`$ is not transversal to $`g^1(c)`$ at $`z`$ if and only if
$$d_zg(T_zE_z)=0\sigma (\stackrel{}{g}(z),T_zE_z)=0\stackrel{}{g}(z)(T_zE_z)^{\mathrm{}}=T_zE_z.$$
$`\mathrm{}`$
If all points of some level $`g^1(c)`$ satisfy conditions of Lemma II.12, then $`g^1(c)`$ is a (2n-1)-dimensional manifold foliated by $`(n1)`$-dimensional submanifolds $`E_zg^1(c)`$. Note that $`\stackrel{}{g}(z)=\mathrm{ker}\sigma |_{T_zg^1(c)}`$, hence $`\mathrm{\Sigma }_z^g\stackrel{def}{=}T_zg^1(c)/\stackrel{}{g}(z)`$ is a $`2(n1)`$-dimensional symplectic space and $`\mathrm{\Delta }_z^g\stackrel{def}{=}T_z\left(E_zg^1(c)\right)`$ is a Lagrangian subspace in $`L_z^g`$, i.e. $`\mathrm{\Delta }_z^gL(\mathrm{\Sigma }_z^g)`$.
The submanifold $`g^1(c)`$ is invariant for the flow $`e^{t\stackrel{}{h}}`$. Moreover, $`e_{}^{t\stackrel{}{h}}\stackrel{}{g}=\stackrel{}{g}`$. Hence $`e_{}^{t\stackrel{}{h}}`$ induces a symplectic transformation $`e_{}^{t\stackrel{}{h}}:\mathrm{\Sigma }_z^g\mathrm{\Sigma }_{e^{t\stackrel{}{h}}(z)}^g`$. Set $`J_z^g(t)=e_{}^{t\stackrel{}{h}}\mathrm{\Delta }_{e^{t\stackrel{}{h}}(z)}^g`$. The curve $`tJ_z^g(t)`$ in the Lagrange Grassmannian $`L(\mathrm{\Sigma }_z^g)`$ is called a reduced Jacobi curve for the Hamiltonian field $`\stackrel{}{h}`$ at $`zN`$.
The reduced Jacobi curve can be easily reconstructed from the Jacobi curve $`J_z(t)=e_{}^{t\stackrel{}{h}}\left(T_{e^{t\stackrel{}{h}}(z)}E_{e^{t\stackrel{}{h}}(z)}\right)L(T_zN)`$ and vector $`\stackrel{}{g}(z)`$. An elementary calculation shows that
$$J_z^g(t)=J_z(t)\stackrel{}{g}(z)^{\mathrm{}}+\stackrel{}{g}(z).$$
Now we can temporary forget the symplectic manifold and Hamiltonians and formulate everything in terms of the curves in the Lagrange Grassmannian. So let $`\mathrm{\Lambda }()`$ be a smooth curve in the Lagrange Grassmannian $`L(\mathrm{\Sigma })`$ and $`\gamma `$ a one-dimensional subspace in $`\mathrm{\Sigma }`$. We set $`\mathrm{\Lambda }^\gamma (t)=\mathrm{\Lambda }(t)\gamma ^{\mathrm{}}+\gamma `$, a Lagrange subspace in the symplectic space $`\gamma ^{\mathrm{}}/\gamma `$. If $`\gamma \mathrm{\Lambda }(t)`$, then $`\mathrm{\Lambda }^\gamma ()`$ is smooth and $`\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}^\gamma (t)=\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}(t)|_{\mathrm{\Lambda }(t)\gamma ^{\mathrm{}}}`$ as it easily follows from the definitions. In particular, monotonicity of $`\mathrm{\Lambda }()`$ implies monotonicity of $`\mathrm{\Lambda }^\gamma ()`$; if $`\mathrm{\Lambda }()`$ is regular and monotone, then $`\mathrm{\Lambda }^\gamma ()`$ is also regular and monotone. The curvatures and the Maslov indices of $`\mathrm{\Lambda }()`$ and $`\mathrm{\Lambda }^\gamma ()`$ are related in a more complicated way. The following result is proved in .
###### Theorem II.2
Let $`\mathrm{\Lambda }(t),t[t_0,t_1]`$ be a smooth monotone increasing curve in $`L(\mathrm{\Sigma })`$ and $`\gamma `$ a one-dimensional subspace of $`\mathrm{\Sigma }`$ such that $`\gamma \mathrm{\Lambda }(t),t[t_0,t_1]`$. Let $`\mathrm{\Pi }L(\mathrm{\Sigma }),\gamma \mathrm{\Pi },\mathrm{\Lambda }(t_0)\mathrm{\Pi }=\mathrm{\Lambda }(t_1)\mathrm{\Pi }=0`$. Then
* $`\mu _\mathrm{\Pi }(\mathrm{\Lambda }())\mu _{\mathrm{\Pi }^\gamma }(\mathrm{\Lambda }^\gamma ())\mu _\mathrm{\Pi }(\mathrm{\Lambda }())+1.`$
* If $`\mathrm{\Lambda }()`$ is regular, then $`r_{\mathrm{\Lambda }^\gamma }(t)r_\mathrm{\Lambda }(t)|_{\mathrm{\Lambda }(t)\gamma ^{\mathrm{}}}`$ and
$`\mathrm{rank}\left(r_{\mathrm{\Lambda }^\gamma }(t)r_\mathrm{\Lambda }(t)|_{\mathrm{\Lambda }(t)\gamma ^{\mathrm{}}}\right)1.`$
The inequality $`r_{\mathrm{\Lambda }^\gamma }(t)r_\mathrm{\Lambda }(t)|_{\mathrm{\Lambda }(t)\gamma ^{\mathrm{}}}`$ turns into the equality if $`\gamma \mathrm{\Lambda }^{}(t),t`$. Then $`\gamma \mathrm{ker}\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}^{}(t)`$. According to Proposition II.9, to $`\gamma `$ there corresponds a one-dimensional subspace in the kernel of $`r_\mathrm{\Lambda }(t)`$; in particular, $`r_\mathrm{\Lambda }(t)`$ is degenerate.
Return to the Jacobi curves $`J_z(t)`$ of a Hamiltonian field $`\stackrel{}{h}`$. There always exists at least one first integral: the Hamiltonian $`h`$ itself. In general, $`\stackrel{}{h}(z)J_z^{}(0)`$ and the reduction procedure has a nontrivial influence on the curvature (see for explicit expressions). Still, there is an important class of Hamiltonians and Lagrange foliations for which the relation $`\stackrel{}{h}(z)J_z^{}(0)`$ holds $`z`$. These are homogeneous on fibers Hamiltonians on cotangent bundles. In this case the generating homotheties of the fibers Euler vector field belongs to the kernel of the curvature form.
### 21 Hyperbolicity
Definition. We say that a Hamiltonian function $`h`$ on the symplectic manifold $`N`$ is regular with respect to the Lagrange foliation $`E`$ if the functions $`h|_{E_z}`$ have nondegenerate second derivatives at $`z,zN`$ (second derivative is well-defined due to the canonical affine structure on $`E_z`$). We say that $`h`$ is monotone with respect to $`E`$ if $`h|_{E_z}`$ is a convex or concave function $`zN`$.
Typical examples of regular monotone Hamiltonians on the cotangent bundles are energy functions of natural mechanical systems. Such a function is the sum of the kinetic energy whose Hamiltonian system generates the Riemannian geodesic flow and a โpotentialโ that is a constant on the fibers function. Proposition II.8 implies that Jacobi curves associated to the regular monotone Hamiltonians are also regular and monotone. Weโll show that negativity of the curvature operators of such a Hamiltonian implies the hyperbolic behavior of the Hamiltonian flow. This is a natural extension of the classical result about Riemannian geodesic flows.
Main tool is the structural equation derived in Section 13. First weโll show that this equation is well coordinated with the symplectic structure. Let $`\mathrm{\Lambda }(t),t,`$ be a regular curve in $`L(\mathrm{\Sigma })`$ and $`\mathrm{\Sigma }=\mathrm{\Lambda }(t)\mathrm{\Lambda }^{}(t)`$ the correspondent canonical splitting. Consider the structural equation
$$\ddot{e}(t)+R_\mathrm{\Lambda }(t)e(t)=0,\mathrm{where}e(t)\mathrm{\Lambda }(t),\dot{e}(t)\mathrm{\Lambda }^{}(t),$$
$`(23)`$
(see Corollary II.1).
###### Lemma II.13
The mapping $`e(0)\dot{e}(0)e(t)\dot{e}(t)`$, where $`e()`$ and $`\dot{e}()`$ satisfies (23), is a symplectic transformation of $`\mathrm{\Sigma }`$.
Proof. We have to check that $`\sigma (e_1(t),e_2(t)),\sigma (\dot{e}_1(t),\dot{e}_2(t)),\sigma (e_1(t),\dot{e}_2(t))`$ do not depend on $`t`$ as soon as $`e_i(t),\dot{e}_i(t),i=1,2`$, satisfy (23). First two quantities vanish since $`\mathrm{\Lambda }(t)`$ and $`\mathrm{\Lambda }^{}(t)`$ are Lagrangian subspaces. The derivative of the third quantity vanishes as well since $`\ddot{e}_i(t)\mathrm{\Lambda }(t).\mathrm{}`$
Let $`h`$ be a regular monotone Hamiltonian on the symplectic manifold $`N`$ equipped with a Lagrange foliation $`E`$. As before, we denote by $`J_z(t)`$ the Jacobi curves of $`\stackrel{}{h}`$ and by $`J_z^h(t)`$ the reduced to the level of $`h`$ Jacobi curves (see previous Section). Let $`R(z)=R_{J_z}(0)`$ and $`R^h(z)=R_{J_z^h}(0)`$ be the curvature operators of $`J_z()`$ and $`J_z^h()`$ correspondently. We say that the Hamiltonian field $`\stackrel{}{h}`$ has a negative curvature at $`z`$ with respect to $`E`$ if all eigenvalues of $`R(z)`$ are negative. We say that $`\stackrel{}{h}`$ has a negative reduced curvature at $`z`$ if all eigenvalues of $`R_z^h`$ are negative.
###### Proposition II.10
Let $`z_0N,z_t=e^{t\stackrel{}{h}}(z)`$. Assume that that $`\overline{\{z_t:t\}}`$ is a compact subset of $`N`$ and that $`N`$ is endowed with a Riemannian structure. If $`\stackrel{}{h}`$ has a negative curvature at any $`z\overline{\{z_t:t\}}`$, then there exists a constant $`\alpha >0`$ and a splitting $`T_{z_t}N=\mathrm{\Delta }_{z_t}^+\mathrm{\Delta }_{z_t}^{}`$, where $`\mathrm{\Delta }_{z_t}^\pm `$ are Lagrangian subspaces of $`T_{z_t}N`$ such that $`e_{}^{\tau \stackrel{}{h}}(\mathrm{\Delta }_{z_t}^\pm )=\mathrm{\Delta }_{z_{t+\tau }}^\pm t,\tau `$ and
$$e_{}^{\pm \tau \stackrel{}{h}}\zeta _\pm e^{\alpha \tau }\zeta _\pm \zeta \mathrm{\Delta }_{z_t}^\pm ,\tau 0,t.$$
$`(24)`$
Similarly, if $`\stackrel{}{h}`$ has a negative reduced curvature at any $`z\overline{\{z_t:t\}}`$, then there exists a splitting $`T_{z_t}(h^1(c)/h(z_t))=\widehat{\mathrm{\Delta }}_{z_t}^+\widehat{\mathrm{\Delta }}_{z_t}^{}`$, where $`c=h(z_0)`$ and $`\widehat{\mathrm{\Delta }}_{z_t}^\pm `$ are Lagrangian subspaces of $`T_{z_t}(h^1(c)/h(z_t))`$ such that $`e_{}^{\tau \stackrel{}{h}}(\widehat{\mathrm{\Delta }}_{z_t}^\pm )=\widehat{\mathrm{\Delta }}_{z_{t+\tau }}^\pm t,\tau `$ and $`e_{}^{\pm \tau \stackrel{}{h}}\zeta _\pm e^{\alpha \tau }\zeta _\pm \zeta \widehat{\mathrm{\Delta }}_{z_t}^\pm ,\tau 0,t.`$
Proof. Obviously, the desired properties of $`\mathrm{\Delta }_{z_t}^\pm `$ and $`\widehat{\mathrm{\Delta }}_{z_t}^\pm `$ do not depend on the choice of the Riemannian structure on $`N`$. Weโll introduce a special Riemannian structure determined by $`h`$. The Riemannian structure is a smooth family of inner products $`,_z`$ on $`T_zN`$, $`zN`$. We have $`T_zN=J_z(0)J_z^{}(0)`$, where $`J_z(0)=T_zE_z`$. Replacing $`h`$ with $`h`$ if necessary we may assume that $`h|_{E_z}`$ is a strongly convex function. First we define $`,_z|_{J_z(0)}`$ to be equal to the second derivative of $`h|_{E_z}`$. Symplectic form $`\sigma `$ induces a nondegenerate pairing of $`J_z(0)`$ and $`J_z^{}(0)`$. In particular, for any $`\zeta J_z(0)`$ there exists a unique $`\zeta ^{}J_z^{}(0)`$ such that $`\sigma (\zeta ^{},)|_{J_z(0)}=\zeta ,_z|_{J_z(0)}`$. There exists a unique extension of the inner product $`,_z`$ from $`J_z(0)`$ to the whole $`T_zN`$ with the following properties:
* $`J_z^{}(0)`$ is orthogonal to $`J_z(0)`$ with respect to $`,_z`$;
* $`\zeta _1,\zeta _2_z=\zeta _1^{},\zeta _2^{}_z\zeta _1,\zeta _2J_z(0)`$.
Weโll need the following classical fact from Hyperbolic Dynamics (see, for instance, \[12, Sec. 17.6\]).
###### Lemma II.14
Let $`A(t),t`$, be a bounded family of symmetric $`n\times n`$-matrices whose eigenvalues are all negative and uniformly separated from 0. Let $`\mathrm{\Gamma }(t,\tau )`$ be the fundamental matrix of the $`2n`$-dimensional linear system $`\dot{x}=y`$, $`\dot{y}=A(t)x`$, where $`x,y^n`$, i.e.
$$\frac{}{t}\mathrm{\Gamma }(t,\tau )=\left(\begin{array}{cc}0& I\\ A& 0\end{array}\right)\mathrm{\Gamma }(t,\tau ),\mathrm{\Gamma }(\tau ,\tau )=\left(\begin{array}{cc}I& 0\\ 0& I\end{array}\right).$$
$`(25)`$
Then there exist closed conic neighborhoods $`C_\mathrm{\Gamma }^+,C_\mathrm{\Gamma }^{}`$, where $`C_\mathrm{\Gamma }^+C_\mathrm{\Gamma }^{}=0`$, of some $`n`$-dimensional subspaces of $`^{2n}`$ and a constant $`\alpha >0`$ such that
$$\mathrm{\Gamma }(t,\tau )C_\mathrm{\Gamma }^+C_\mathrm{\Gamma }^+,|\mathrm{\Gamma }(t,\tau )\xi _+|e^{\alpha (\tau t)}|\xi _+|,\xi _+C_\mathrm{\Gamma }^+,t\tau ,$$
and
$$\mathrm{\Gamma }(t,\tau )C_\mathrm{\Gamma }^{}C_\mathrm{\Gamma }^{},|\mathrm{\Gamma }(t,\tau )\xi _{}|e^{\alpha (t\tau )}|\xi _{}|,\xi _{}C_\mathrm{\Gamma }^{},t\tau .$$
The constant $`\alpha `$ depends only on upper and lower bounds of the eigenvalues of $`A(t).\mathrm{}`$
###### Corollary II.6
Let $`C_\mathrm{\Gamma }^\pm `$ be the cones described in Lemma II.14; then $`\mathrm{\Gamma }(0,\pm t)C_\mathrm{\Gamma }^\pm \mathrm{\Gamma }(0;\pm \tau )C_\mathrm{\Gamma }^\pm `$ for any $`t\tau 0`$ and the subsets $`K_\mathrm{\Gamma }^\pm =\underset{t0}{}\mathrm{\Gamma }(0,t)C_\mathrm{\Gamma }^\pm `$ are Lagrangian subspaces of $`^n\times ^n`$ equipped with the standard symplectic structure.
Proof. The relations $`\mathrm{\Gamma }(\tau ,t)C_\mathrm{\Gamma }^+C_\mathrm{\Gamma }^+`$ and $`\mathrm{\Gamma }(\tau ,t)C_\mathrm{\Gamma }^{}C_\mathrm{\Gamma }^{}`$ imply:
$$\mathrm{\Gamma }(0,\pm t)C_\mathrm{\Gamma }^\pm =\mathrm{\Gamma }(0,\pm \tau )\mathrm{\Gamma }(\pm \tau ,\pm t)C_\mathrm{\Gamma }^\pm \mathrm{\Gamma }(0,\pm \tau )C_\mathrm{\Gamma }^\pm .$$
In what follows weโll study $`K_\mathrm{\Gamma }^+`$; the same arguments work for $`K_\mathrm{\Gamma }^{}`$. Take vectors $`\zeta ,\zeta ^{}K_\mathrm{\Gamma }^+`$; then $`\zeta =\mathrm{\Gamma }(0,t)\zeta _t`$ and $`\zeta ^{}=\mathrm{\Gamma }(0,t)\zeta _t^{}`$ for any $`t0`$ and some $`\zeta _t,\zeta _t^{}C_\mathrm{\Gamma }^+`$. Then, according to Lemma II.14, $`|\zeta _t|e^{\alpha t}|\zeta |,|\zeta _t^{}|e^{\alpha t}|\zeta ^{}|`$, i.e. $`\zeta _t`$ and $`\zeta _t^{}`$ tend to 0 as $`t+\mathrm{}`$. On the other hand,
$$\sigma (\zeta ,\zeta ^{})=\sigma (\mathrm{\Gamma }(0,t)\zeta _t,\mathrm{\Gamma }(0,t)\zeta _t^{})=\sigma (\zeta _t,\zeta _t^{})t0$$
since $`\mathrm{\Gamma }(0,t)`$ is a symplectic matrix. Hence $`\sigma (\zeta ,\zeta ^{})=0`$.
We have shown that $`K_\mathrm{\Gamma }^+`$ is an isotropic subset of $`^n\times ^n`$. On the other hand, $`K_\mathrm{\Gamma }^+`$ contains an $`n`$-dimensional subspace since $`C_\mathrm{\Gamma }^+`$ contains one and $`\mathrm{\Gamma }(0,t)`$ are invertible linear transformations. Isotropic $`n`$-dimensional subspace is equal to its skew-orthogonal complement, therefore $`K_\mathrm{\Gamma }^+`$ is a Lagrangian subspace. $`\mathrm{}`$
Take now a regular monotone curve $`\mathrm{\Lambda }(t),t`$ in the Lagrange Grassmannian $`L(\mathrm{\Sigma })`$. We may assume that $`\mathrm{\Lambda }()`$ is monotone increasing, i.e. $`\dot{\mathrm{\Lambda }}(t)>0`$. Recall that $`\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}(t)(e(t))=\sigma (e(t),\dot{e}(t))`$, where $`e()`$ is an arbitrary smooth curve in $`\mathrm{\Sigma }`$ such that $`e(\tau )\mathrm{\Lambda }(\tau ),\tau `$. Differentiation of the identity $`\sigma (e_1(\tau ),e_2(\tau ))=0`$ implies: $`\sigma (e_1(t),\dot{e}_2(t))=\sigma (\dot{e}_1(t),e_2(t))=\sigma (e_2(t),\dot{e}_1(t))`$ if $`e_i(\tau )\mathrm{\Lambda }(\tau )`$, $`\tau `$, $`i=1,2`$. Hence the Euclidean structure $`,_{\dot{\mathrm{\Lambda }}(t)}`$ defined by the quadratic form $`\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}(t)`$ reads: $`e_1(t),e_2(t)_{\dot{\mathrm{\Lambda }}(t)}=\sigma (e_1(t),\dot{e}_2(t))`$.
Take a basis $`e_1(0),\mathrm{},e_n(0)`$ of $`\mathrm{\Lambda }(0)`$ such that the form $`\underset{ยฏ}{\overset{\dot{}}{\mathrm{\Lambda }}}(t)`$ has the unit matrix in this basis, i.e. $`\sigma (e_i(0),\dot{e}_j(0))=\delta _{ij}`$. In fact, vectors $`\dot{e}_j(0)`$ are defined modulo $`\mathrm{\Lambda }(0)`$; we can normalize them assuming that $`\dot{e}_i(0)\mathrm{\Lambda }^{}(0),i=1,\mathrm{},n`$. Then $`e_1(0),\mathrm{},e_n(0),\dot{e}_1(0),\mathrm{},\dot{e}_n(0)`$ is a Darboux basis of $`\mathrm{\Sigma }`$. Fix coordinates in $`\mathrm{\Sigma }`$ using this basis: $`\mathrm{\Sigma }=^n\times ^n`$, where $`\left(\begin{array}{c}x\\ y\end{array}\right)^n\times ^n`$ is identified with $`\underset{j=1}{\overset{n}{}}\left(x^je_j(0)+y^j\dot{e}_j(0)\right)\mathrm{\Sigma },`$ $`x=(x^1,\mathrm{},x^n)^{}`$, $`y=(y^1,\mathrm{},y^n)^{}.`$
We claim that there exists a smooth family $`A(t),t,`$ of symmetric $`n\times n`$ matrices such that $`A(t)`$ has the same eigenvalues as $`R_\mathrm{\Lambda }(t)`$ and
$$\mathrm{\Lambda }(t)=\mathrm{\Gamma }(0,t)\left(\begin{array}{c}^n\\ 0\end{array}\right),\mathrm{\Lambda }^{}(t)=\mathrm{\Gamma }(0,t)\left(\begin{array}{c}0\\ ^n\end{array}\right),t$$
in the fixed coordinates, where $`\mathrm{\Gamma }(t,\tau )`$ satisfies (25). Indeed, let $`e_i(t),i=1,\mathrm{},n,`$ be solutions to the structural equations (23). Then
$$\mathrm{\Lambda }(t)=span\{e_1(t),\mathrm{},e_n(t)\},\mathrm{\Lambda }^{}(t)=span\{\dot{e}_1(t),\mathrm{},\dot{e}_n(t)\}.$$
Moreover, $`\ddot{e}_i(t)=\underset{i=1}{\overset{n}{}}a_{ij}(t)e_j(t)`$, where $`A(t)=\{a_{ij}(t)\}_{i,j=1}^n`$ is the matrix of the operator $`R_\mathrm{\Lambda }(t)`$ in the โmovingโ basis $`e_1(t),\mathrm{},e_n(t)`$. Lemma I.13 implies that $`e_i(t),e_j(t)_{\dot{\mathrm{\Lambda }}(t)}=\sigma (e_i(t),\dot{e}_j(t))=\delta _{ij}`$. In other words, the Euclidean structure $`,_{\dot{\mathrm{\Lambda }}(t)}`$ has unit matrix in the basis $`e_1(t),\mathrm{},e_n(t)`$. Operator $`R_\mathrm{\Lambda }(t)`$ is self-adjoint for the Euclidean structure $`,_{\dot{\mathrm{\Lambda }}(t)}`$ (see Propositon II.9). Hence matrix $`A(t)`$ is symmetric.
Let $`e_i(t)=\left(\begin{array}{c}x_i\left(t\right)\\ y_i\left(t\right)\end{array}\right)^n\times ^n`$ in the fixed coordinates. Make up $`n\times n`$-matrices $`X(t)=(x_1(t),\mathrm{},x_n(t))`$, $`Y(t)=(y_1(t),\mathrm{},y_n(t))`$ and a $`2n\times 2n`$-matrix $`\left(\begin{array}{cc}X\left(t\right)& \dot{X}\left(t\right)\\ Y\left(t\right)& \dot{Y}\left(t\right)\end{array}\right).`$ We have
$$\frac{d}{dt}\left(\begin{array}{cc}X& \dot{X}\\ Y& \dot{Y}\end{array}\right)(t)=\left(\begin{array}{cc}X& \dot{X}\\ Y& \dot{Y}\end{array}\right)(t)\left(\begin{array}{cc}0& A(t)\\ I& 0\end{array}\right),\left(\begin{array}{cc}X& \dot{X}\\ Y& \dot{Y}\end{array}\right)(0)=\left(\begin{array}{cc}I& 0\\ 0& I\end{array}\right).$$
Hence $`\left(\begin{array}{cc}X& \dot{X}\\ Y& \dot{Y}\end{array}\right)(t)=\mathrm{\Gamma }(t,0)^1=\mathrm{\Gamma }(0,t)`$.
Let now $`\mathrm{\Lambda }()`$ be the Jacobi curve, $`\mathrm{\Lambda }(t)=J_{z_0}(t)`$. Set $`\xi _i(z_t)=e_{}^{t\stackrel{}{h}}e_i(t)`$, $`\eta _i(z_t)=e_{}^{t\stackrel{}{h}}\dot{e}_i(t)`$; then
$$\xi _1(z_t),\mathrm{},\xi _n(z_t),\eta _1(z_t),\mathrm{},\eta _n(z_t)$$
$`(26)`$
is a Darboux basis of $`T_{z_t}N`$, where $`J_{z_t}(0)=span\{\xi _1(z_t),\mathrm{},\xi _n(z_t)\}`$, $`J_{z_t}^{}(0)=span\{\eta _1(z_t),\mathrm{},\eta _n(z_t)\}`$. Moreover, the basis (26) is orthonormal for the inner product $`,_{z_t}`$ on $`T_{z_t}N`$.
The intrinsic nature of the structural equation implies the translation invariance of the construction of the frame (26): if we would start from $`z_s`$ instead of $`z_0`$ and put $`\mathrm{\Lambda }(t)=J_{z_s}(t)`$, $`e_i(0)=\xi _i(z_s)`$, $`\dot{e}_i(0)=\eta _i(z_s)`$ for some $`s`$, then we would obtain $`e_{}^{t\stackrel{}{h}}e_i(t)=\xi _i(z_{s+t})`$, $`e_{}^{t\stackrel{}{h}}\dot{e}_i(t)=\eta _i(z_{s+t})`$.
The frame (26) gives us fixed orthonormal Darboux coordinates in $`T_{z_s}N`$ for $`s`$ and the correspondent symplectic $`2n\times 2n`$-matrices $`\mathrm{\Gamma }_{z_s}(\tau ,t)`$. We have: $`\mathrm{\Gamma }_{z_s}(\tau ,t)==\mathrm{\Gamma }_{z_0}(s+\tau ,s+t)`$; indeed, $`\mathrm{\Gamma }_{z_s}(\tau ,t)\left(\begin{array}{c}x\\ y\end{array}\right)`$ is the coordinate presentation of the vector
$$e_{}^{(\tau t)\stackrel{}{h}}\left(\underset{i}{}x^i\xi ^i(z_{s+t})+y^i\eta _i(z_{s+t})\right)$$
in the basis $`\xi _i(z_{s+\tau }),\eta _i(z_{s+\tau })`$. In particular,
$$\left|\mathrm{\Gamma }_{z_s}(0,t)\left(\begin{array}{c}x\\ y\end{array}\right)\right|=e_{}^{t\stackrel{}{h}}\left(\underset{i}{}x^i\xi ^i(z_{s+t})+y^i\eta _i(z_{s+t})\right)_{z_s}.$$
$`(27)`$
Recall that $`\xi _1(z_\tau ),\mathrm{},\xi _n(z_\tau ),\eta _1(z_\tau ),\mathrm{},\eta _n(z_\tau )`$ is an orthonormal frame for the scalar product $`,_{z_\tau }`$ and $`\zeta _{z_\tau }=\sqrt{\zeta ,\zeta }_{z_\tau }`$.
We introduce the notation :
$$W_{z_s}=\{\underset{i}{}x^i\xi ^i(z_s)+y^i\eta _i(z_s):\left(\begin{array}{c}x\\ y\end{array}\right)W\},$$
for any $`W^n\times ^n`$. Let $`C_{\mathrm{\Gamma }_{z_0}}^\pm `$ be the cones from Lemma II.14. Then
$$e_{}^{\tau \stackrel{}{h}}\mathrm{\Gamma }_{z_s}(0,t)C_{\mathrm{\Gamma }_{z_0}}^\pm _{z_{s\tau }}=\mathrm{\Gamma }_{z_{s\tau }}(0,t+\tau )C_{\mathrm{\Gamma }_{z_0}}^\pm _{z_{s\tau }},t,\tau ,s.$$
$`(28)`$
Now set $`K_{\mathrm{\Gamma }_{z_s}}^+=\underset{t0}{}C_{\mathrm{\Gamma }_{z_0}}^+`$, $`K_{\mathrm{\Gamma }_{z_s}}^{}=\underset{t0}{}C_{\mathrm{\Gamma }_{z_0}}^{}`$ and $`\mathrm{\Delta }_{z_s}^\pm =K_{\mathrm{\Gamma }_{z_s}}^{}_{z_s}`$. Corollary II.6 implies that $`\mathrm{\Delta }_{z_s}^\pm `$ are Lagrangian subspaces of $`T_{z_s}N`$. Moreover, it follows from (28) that $`e_{}^{t\stackrel{}{h}}\mathrm{\Delta }_{z_s}^\pm =\mathrm{\Delta }_{z_{s+t}}^\pm `$, while (28) and (27) imply inequalities (24).
This finishes the proof of the part of Proposition II.10 which concerns Jacobi curves $`J_z(t)`$. We leave to the reader a simple adaptation of this proof to the case of reduced Jacobi curves $`J_z^h(t).\mathrm{}`$
Remark. Constant $`\alpha `$ depends, of course, on the Riemannian structure on $`N`$. In the case of the special Riemannian structure defined at the beginning of the proof of Proposition II.10 this constant depends only on the upper and lower bounds for the eigenvalues of the curvature operators and reduced curvature operators correspondently (see Lemma II.14 and further arguments).
Let $`e^{tX},t`$ be the flow generated by the the vector field $`X`$ on a manifold $`M`$. Recall that a compact invariant subset $`WM`$ of the flow $`e^{tX}`$ is called a hyperbolic set if there exists a Riemannian structure in a neighborhood of $`W`$, a positive constant $`\alpha `$, and a splitting $`T_zM=E_z^+E_z^{}X(z),zW`$, such that $`X(z)0,e_{}^{tX}E_z^\pm =E_{e^{tX}(z)}^\pm `$, and $`e_{}^{\pm tX}\zeta ^\pm e^{\alpha t}\zeta ^\pm ,t0,\zeta ^\pm E_z^\pm `$. Just the fact some invariant set is hyperbolic implies a rather detailed information about asymptotic behavior of the flow in a neighborhood of this set (see for the introduction to Hyperbolic Dynamics). The flow $`e^{tX}`$ is called an Anosov flow if the entire manifold $`M`$ is a hyperbolic set.
The following result is an immediate corollary of Proposition II.10 and the above remark.
###### Theorem II.3
Let $`h`$ be a regular monotone Hamiltonian on $`N`$, $`c`$, $`Wh^1(c)`$ a compact invariant set of the flow $`e^{t\stackrel{}{h}},t`$, and $`d_zh0,zW`$. If $`\stackrel{}{h}`$ has a negative reduced curvature at every point of $`W`$, then $`W`$ is a hyperbolic set of the flow $`e^{t\stackrel{}{h}}|_{h^1(c)}.\mathrm{}`$
This theorem generalizes a classical result about geodesic flows on compact Riemannian manifolds with negative sectional curvatures. Indeed, if $`N`$ is the cotangent bundle of a Riemannian a Riemannian manifold and $`e^{t\stackrel{}{h}}`$ is the geodesic flow, then negativity of the reduced curvature of $`\stackrel{}{h}`$ means simply negativity of the sectional Riemannian curvature. In this case, the Hamiltonian $`h`$ is homogeneous on the fibers of the cotangent bundle and the restrictions $`e^{t\stackrel{}{h}}|_{h^1(c)}`$ are equivalent for all $`c>0`$.
The situation changes if $`h`$ is the energy function of a general natural mechanical system on the Riemannian manifold. In this case, the flow and the reduced curvature depend on the energy level. Still, negativity of the sectional curvature implies negativity of the reduced curvature at $`h^1(c)`$ for all sufficiently big $`c`$. In particular, $`e^{t\stackrel{}{h}}|_{h^1(c)}`$ is an Anosov flow for any sufficiently big $`c`$; see for the explicit expression of the reduced curvature in this case.
Theorem II.3 concerns only the reduced curvature while the next result deals with the (not reduced) curvature of $`\stackrel{}{h}`$.
###### Theorem II.4
Let $`h`$ be a regular monotone Hamiltonian and $`W`$ a compact invariant set of the flow $`e^{t\stackrel{}{h}}`$. If $`\stackrel{}{h}`$ has a negative curvature at any point of $`W`$, then $`W`$ is a finite set and each point of $`W`$ is a hyperbolic equilibrium of the field $`\stackrel{}{h}`$.
Proof. Let $`zW`$; the trajectory $`z_t=e^{t\stackrel{}{h}}(z),t`$, satisfies conditions of Proposition II.10. Take the correspondent splitting $`T_{z_t}N=\mathrm{\Delta }_{z_t}^+\mathrm{\Delta }_{z_t}^{}`$. In particular, $`\stackrel{}{h}(z_t)=\stackrel{}{h}^+(z_t)+\stackrel{}{h}^{}(z_t)`$, where $`\stackrel{}{h}^\pm (z_t)\mathrm{\Delta }_{z_t}^\pm `$.
We have $`e_{}^{\tau \stackrel{}{h}}\stackrel{}{h}(z_t)=\stackrel{}{h}(z_{t+\tau })`$. Hence
$$\stackrel{}{h}(z_{t+\tau })=e_{}^{\tau \stackrel{}{h}}\stackrel{}{h}(z_t)e_{}^{\tau \stackrel{}{h}}\stackrel{}{h}^+(z_t)e_{}^{\tau \stackrel{}{h}}\stackrel{}{h}^{}(z_t)$$
$$e^{\alpha \tau }\stackrel{}{h}^+(z_t)e^{\alpha \tau }\stackrel{}{h}^{}(z_t),\tau 0.$$
Compactness of $`\overline{\{z_t:t\}}`$ implies that $`\stackrel{}{h}^+(z_t)`$ is uniformly bounded; hence $`\stackrel{}{h}^+(z_t)=0`$. Similarly, $`\stackrel{}{h}(z_{t\tau }e^{\alpha \tau }\stackrel{}{h}^{}(z_t)e^{\alpha \tau }\stackrel{}{h}^+(z_t)`$ that implies the equality $`\stackrel{}{h}^{}(z_t)=0`$. Finally, $`\stackrel{}{h}(z_t)=0`$. In other words, $`z_tz`$ is an equilibrium of $`\stackrel{}{h}`$ and $`T_zN=\mathrm{\Delta }_z^+\mathrm{\Delta }_z^{}`$ is the splitting of $`T_zN`$ into the repelling and attracting invariant subspaces for the linearization of the flow $`e^{t\stackrel{}{h}}`$ at $`z`$. Hence $`z`$ is a hyperbolic equilibrium; in particular, $`z`$ is an isolated equilibrium of $`\stackrel{}{h}.\mathrm{}`$
We say that a subset of a finite dimensional manifold is bounded if it has a compact closure.
###### Corollary II.7
Assume that $`h`$ is a regular monotone Hamiltonian and $`\stackrel{}{h}`$ has everywhere negative curvature. Then any bounded semi-trajectory of the system $`\dot{z}=\stackrel{}{h}(z)`$ converges to an equilibrium with the exponential rate while another semi-trajectory of the same trajectory must be unbounded. $`\mathrm{}`$
Typical Hamiltonians which satisfy conditions of Corollary II.7 are energy functions of natural mechanical systems in $`^n`$ with a strongly concave potential energy. Indeed, in this case, the second derivative of the potential energy is equal to the matrix of the curvature operator in the standard Cartesian coordinates (see Sec. 15).
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# A Garside-theoretic approach to the reducibility problem in braid groups
## 1. Introduction
Let $`D_n=\{z:|z|n+1\}\{1,\mathrm{},n\}`$, the $`n`$-punctured disk in the complex plane with punctures lying on the real axis. The $`n`$-braid group $`B_n`$ acts on the set of curve systems in $`D_n`$. For an $`n`$-braid $`\alpha `$ and a curve system $`๐`$ in $`D_n`$, let $`\alpha ๐`$ denote the action of $`\alpha `$ on $`๐`$. An $`n`$-braid $`\alpha `$ is said to be *reducible* if $`\alpha ๐=๐`$ for some essential curve system $`๐`$ in $`D_n`$, called a *reduction system* of $`\alpha `$. In this paper, we are interested in the *reducibility problem*: given a braid, decide whether it is reducible or not and find a reduction system if it is reducible.
### 1.1. Motivation and some of previous works
The Nielsen-Thurston classification theorem \[Thu88\] states that an irreducible automorphism of an orientable surface with negative euler characteristic is either periodic or pseudo-Anosov up to isotopy. Recall that an orientation preserving self-diffeomorphism $`f`$ of a surface $`S`$ is said to be
* *periodic* if $`f^k`$ is isotopic to the identity for some $`k0`$;
* *reducible* if there exist pairwise disjoint simple closed curves $`C_1,\mathrm{},C_k`$ in $`S`$, isotopic to neither a point nor a puncture nor a boundary component, such that $`f(๐)`$ is isotopic to $`๐`$, where $`๐=C_1\mathrm{}C_k`$;
* *pseudo-Anosov* if there exist a pair of transverse measured foliations $`(F^s,\mu ^s)`$ and $`(F^u,\mu ^u)`$ and a real $`\lambda >1`$ such that $`f(F^s,\mu ^s)=(F^s,\lambda ^1\mu ^s)`$ and $`f(F^u,\mu ^u)=(F^u,\lambda \mu ^u)`$.
There have been several approaches to the problem of deciding dynamical types of surface automorphisms. Bestvina and Handel \[BH95\] made the train track algorithm that, given a surface automorphism, decides its dynamical type and finds its dynamical structure: a pair of transverse measured foliations for a pseudo-Anosov automorphism; a reduction system for a reducible automorphism. Benardete, Gutiรฉrrez and Nitecki \[BGN95\] solved the reducibility problem in braid groups. (It is known that a periodic $`n`$-braid is conjugate to either $`(\sigma _1\sigma _2\mathrm{}\sigma _{n1})^l`$ or $`(\sigma _1(\sigma _1\sigma _2\mathrm{}\sigma _{n1}))^l`$ for some integer $`l`$ \[Ker19, Eil34, BDM02\]. This implies that $`\alpha `$ is a periodic $`n`$-braid if and only if either $`\alpha ^n`$ or $`\alpha ^{n1}`$ is equal to $`\mathrm{\Delta }^{2m}`$ for some integer $`m`$. Hence, it is easy to decide the periodicity of braids. Therefore, in order to decide the dynamical type of a given braid, it suffices to decide the reducibility.) Humphries \[Hum91\] solved the problem of recognizing split braids.
With the above results, solving the reducibility problem and the problem of recognizing split braids seems at least as hard as solving the conjugacy problem. When using the train track algorithm, one needs to describe a given $`n`$-braid as a graph map of the $`n`$-bouquet, and the length of this description grows exponentially with respect to the length of the braid word on Artin generators. The other two solutions need to use the algorithms solving the conjugacy problem in braid groups.
Another motivation for this work is the close relationship between the reducibility problem and the conjugacy problem. The approach to the conjugacy problem in braid groups can be divided into two steps: solving the reducibility problem and solving the conjugacy problem for irreducible braids. See \[BGG06a, ยง1.4\] for a more precise description of this strategy. The conjugacy problem for periodic braids is easy to solve. There are two different polynomial-time solutions to this case by Birman, Gebhardt and Gonzรกlez-Meneses \[BGG06b\] and by the authors \[LL07b\]. For the conjugacy problem for pseudo-Anosov mapping classes, there are several results. In \[Los93\], Los solved the problem for pseudo-Anosov braids by using combinatorial efficient representatives. Recently, Fehrenbach and Los \[FL07\] proposed an algorithm that finds roots and symmetries of pseudo-Anosov mapping classes together with a new solution to the conjugacy problem. Mazur and Minsky \[MM99, MM00\] showed that, fixing a mapping class group and a finite set of generators for this group, there exists a constant $`K`$ such that if $`\alpha `$ and $`\beta `$ are conjugate pseudo-Anosov mapping classes then there is a conjugating element $`\gamma `$ with $`|\gamma |K(|\alpha |+|\beta |)`$, where $`||`$ denotes the word length. In order to extend the results on irreducible braids to general braids, we need to solve the reducibility problem more efficiently.
For the last ten years, no serious progress has been made in the reducibility problem. On the other hand, recently, there have been several new contributions to Garside-theoretic approach to braid groups, for example \[Deh02, FG03, Geb05, Lee07\]. Exploiting them, we study the characteristics of the conjugacy classes of reducible braids. Our approach uses neither the train track algorithm nor the complete conjugacy algorithm. We hope that our results are useful in obtaining a more efficient solution to the reducibility problem in braid groups.
### 1.2. Our results
Before stating our results, we recall some notions and results from the Garside theory in braid groups.
* Let $`B_n^+`$ be the submonoid of $`B_n`$ generated by $`\sigma _1,\mathrm{},\sigma _{n1}`$. The partial order $`_R`$ on $`B_n^+`$ is defined as follows: for $`P,QB_n^+`$, $`P_RQ`$ if $`Q=SP`$ for some $`SB_n^+`$. The poset $`(B_n^+,_R)`$ is a lattice, i.e., there exist the gcd $`P_RQ`$ and the lcm $`P_RQ`$ of $`P,QB_n^+`$.
* For $`\alpha B_n`$, there are integer-valued invariants $`inf(\alpha )`$ and $`sup(\alpha )`$. Let $`[\alpha ]`$ denote the conjugacy class of $`\alpha B_n`$. The following are conjugacy invariants.
$$\begin{array}{cc}inf_s(\alpha )=\mathrm{max}\{inf(\beta ):\beta [\alpha ]\}\hfill & t_{inf}(\alpha )=lim_m\mathrm{}inf(\alpha ^m)/m\hfill \\ sup_s(\alpha )=\mathrm{min}\{sup(\beta ):\beta [\alpha ]\}\hfill & t_{sup}(\alpha )=lim_m\mathrm{}sup(\alpha ^m)/m\hfill \end{array}$$
* In the conjugacy class $`[\alpha ]`$, there are finite, nonempty, computable subsets, the super summit set $`[\alpha ]^S`$, the ultra summit set $`[\alpha ]^U`$ and the stable super summit set $`[\alpha ]^{St}`$. They depend only on the conjugacy class, and $`[\alpha ]^U,[\alpha ]^{St}[\alpha ]^S`$.
We call an essential curve system (see Definition 3.1) in $`D_n`$ a *standard* curve system if each component is isotopic to a round circle centered at the real axis as in Figure 1.
For an essential curve system $`๐`$ in $`D_n`$, we define the *standardizer* of $`๐`$ as the set
$$\mathrm{St}(๐)=\{PB_n^+:P๐\text{ is standard}\}$$
where $`P๐`$ denotes the left action of the positive braid $`P`$ on the curve system $`๐`$, and then show the following.
Theorem 4.2. For an essential curve system $`๐`$ in $`D_n`$, its standardizer $`\mathrm{St}(๐)`$ is closed under $`_R`$ and $`_R`$, and hence a sublattice of $`B_n^+`$. Therefore $`\mathrm{St}(๐)`$ contains a unique $`_R`$-minimal element.
Theorem 4.9. Let $`\alpha `$ be a reducible $`n`$-braid with a reduction system $`๐`$. Let $`P`$ be the $`_R`$-minimal element of $`\mathrm{St}(๐)`$. Then the following hold.
1. $`inf(\alpha )inf(P\alpha P^1)sup(P\alpha P^1)sup(\alpha )`$.
2. If $`\alpha [\alpha ]^S`$, then $`P\alpha P^1[\alpha ]^S`$.
3. If $`\alpha [\alpha ]^U`$, then $`P\alpha P^1[\alpha ]^U`$.
4. If $`\alpha [\alpha ]^{St}`$, then $`P\alpha P^1[\alpha ]^{St}`$.
Theorem 4.2 is essential in our approach to the reducibility problem, as the closedness under $`_R`$ of $`\{PB_n^+:P\beta P^1[\alpha ]^S\}`$ and $`\{PB_n^+:P\beta P^1[\alpha ]^U\}`$ for $`\beta [\alpha ]^S`$ plays an important role in solving the conjugacy problem \[FG03, Geb05\]. Theorem 4.9 shows that standardizing a reduction system $`๐`$ of a braid by the $`_R`$-minimal element of $`\mathrm{St}(๐)`$ preserves the membership of the super summit set, ultra summit set and stable super summit set.
It is known by Birman, Lubotzky and McCarthy \[BLM83\] and Ivanov \[Iva92\] that a reducible surface automorphism admits a unique *canonical reduction system*. For $`\alpha B_n`$, let $`_{\mathrm{ext}}(\alpha )`$ be the collection of the outermost components of the canonical reduction system of $`\alpha `$. Let $`P`$ be the $`_R`$-minimal element of $`\mathrm{St}(_{\mathrm{ext}}(\alpha ))`$. Since $`_{\mathrm{ext}}(P\alpha P^1)=P_{\mathrm{ext}}(\alpha )`$ is standard, the outermost component of $`D_n_{\mathrm{ext}}(P\alpha P^1)`$ is naturally identified with the $`k`$-punctured disk $`D_k`$ for some $`kn`$. We define the *outermost component* $`\alpha _{\mathrm{ext}}`$ of $`\alpha `$ as the $`k`$-braid obtained by restricting the braid $`P\alpha P^1`$ to the outermost component of $`D_n_{\mathrm{ext}}(P\alpha P^1)`$. See ยง5 for the precise definition. The following is the main result of this paper. (In the statement, $`[\alpha ]_๐^U`$ denotes the ultra summit set of $`\alpha `$ with respect to decycling. See the next section for the precise definition.)
Theorem 7.4. Let $`\alpha `$ be a non-periodic reducible $`n`$-braid.
1. If $`inf_s(\alpha _{\mathrm{ext}})>inf_s(\alpha )`$, then each element of $`[\alpha ]^U`$ has a standard reduction system.
2. If $`sup_s(\alpha _{\mathrm{ext}})<sup_s(\alpha )`$, then each element of $`[\alpha ]_๐^U`$ has a standard reduction system.
3. If $`\alpha `$ is a split braid, then each element of $`[\alpha ]^U[\alpha ]_๐^U`$ has a standard reduction system.
4. If $`\alpha _{\mathrm{ext}}`$ is periodic, then there exists $`1q<n`$ such that each element of $`[\alpha ^q]^U[\alpha ^q]_๐^U`$ has a standard reduction system.
5. If $`t_{inf}(\alpha _{\mathrm{ext}})>t_{inf}(\alpha )`$, then there exists $`1q<n(n1)/2`$ such that each element of $`[\alpha ^q]^U`$ has a standard reduction system.
6. If $`t_{sup}(\alpha _{\mathrm{ext}})<t_{sup}(\alpha )`$, then there exists $`1q<n(n1)/2`$ such that each element of $`[\alpha ^q]_๐^U`$ has a standard reduction system.
Roughly speaking, the first statement of the above theorem says that if the outermost component $`\alpha _{\mathrm{ext}}`$ is simpler than the whole braid $`\alpha `$ up to conjugacy from a Garside-theoretic point of view, then every element of $`[\alpha ]^U`$ has a standard reduction system. In this case, finding a reduction system is as easy as finding one element in the ultra summit set, because it is easy to find a standard reduction system of a given braid if it exists by the results in \[BGN93\]. In ยง7, we present three examples showing that the conditions in Theorem 7.4 cannot be weakened.
In \[BGN95\], Benardete, Gutiรฉrrez and Nitecki showed that *if a braid is reducible, then there exists an element in its super summit set which has a standard reduction system.* (The notion of ultra summit set appeared later than their work, and from their proof we can replace โsuper summit setโ in their statement with โultra summit setโ.) While their result concerns the *existence* of an ultra summit element with a standard reduction system, Theorem 7.4 (i)-(iii) show that, under a certain condition, *every* ultra summit element has a standard reduction system.
We remark that the six types of braids in Theorem 7.4 cover most reducible braids. The braid $`\alpha _{\mathrm{ext}}`$ can be obtained, up to conjugacy, by deleting some strands from $`\alpha `$, hence $`\alpha _{\mathrm{ext}}`$ cannot be more complicated than $`\alpha `$. Indeed, the following inequalities always hold (see Lemma 5.3):
$$\begin{array}{cc}inf_s(\alpha _{\mathrm{ext}})inf_s(\alpha );& sup_s(\alpha _{\mathrm{ext}})sup_s(\alpha );\\ t_{inf}(\alpha _{\mathrm{ext}})t_{inf}(\alpha );& t_{sup}(\alpha _{\mathrm{ext}})t_{sup}(\alpha ).\end{array}$$
Theorem 7.4 shows the characteristics of the braid conjugacy classes for which at least one of the above inequalities is strict.
We briefly explain the idea of proof of Theorem 7.4.
* In ยง6, we show that if $`\alpha `$ is a split braid with the minimal word length in the conjugacy class, then the outermost component $`_{\mathrm{ext}}(\alpha )`$ of the canonical reduction system of $`\alpha `$ is standard. Since a positive braid has the minimal word length in the conjugacy class, we have the following: *if $`P`$ is a positive split braid, then $`_{\mathrm{ext}}(P)`$ is standard.*
* If a braid $`\alpha `$ commutes with a non-periodic reducible braid $`\beta `$, then the canonical reduction system of $`\beta `$ is a reduction system of $`\alpha `$. Combining this with the previous observation, we have the following: *if $`\alpha P=P\alpha `$ for some positive split braid $`P`$, then $`_{\mathrm{ext}}(P)`$ is a standard reduction system of $`\alpha `$.*
* If $`\alpha `$ belongs to the ultra summit set, then there exists a finite sequence $`\alpha =\alpha _0\alpha _1\mathrm{}\alpha _m=\alpha `$ for some $`m1`$, where $`\alpha _{i+1}=A_i\alpha _iA_i^1`$ for some permutation braid $`A_i`$ for $`i=0,\mathrm{},m1`$. If we let $`T=A_{m1}\mathrm{}A_1A_0`$, then $`T\alpha =\alpha T`$. Exploiting the $`_R`$-minimal elements of the standardizers $`\mathrm{St}(_{\mathrm{ext}}(\alpha _i))`$, we show that $`T`$ is a positive split braid if $`inf_s(\alpha _{\mathrm{ext}})>inf_s(\alpha )`$, from which Theorem 7.4 (i) follows. The other statements are proved using this.
### 1.3. Organization
In ยง2, we review the Garside theory in brad groups. In ยง3, we study the normal form of the braids that send a standard curve system to a standard curve system. In ยง4, we prove Theorems 4.2 and 4.9. In ยง5, we study the properties of the outermost component $`\alpha _{\mathrm{ext}}`$ of a non-periodic reducible braid $`\alpha `$. In ยง6, we show that if a split braid has the minimal word length in the conjugacy class, then the outermost component of its canonical reduction system is standard. In ยง7 and ยง8, we prove Theorem 7.4, using the results of the previous sections.
### Acknowledgements
We are most grateful to the anonymous referee of this journal for valuable comments and suggestions on the paper, especially for pointing out that our initial proof of Theorem 4.9 contains a mistake. The proof is corrected as suggested by the referee. We are also very thankful to Won Taek Song for helpful conversations, and Juan Gonzรกlez-Meneses and Bert Wiest for providing Example 7.7. This work was supported by the Korea Science and Engineering Foundation (KOSEF) grant funded by the Korea government (MOST) (No. R01-2007-000-20293-0).
## 2. Garside theory in braid groups
We give necessary definitions and results on Garside theory in braid groups. See \[Gar69, Thu92, EM94, BKL98, DP99, Deh02, FG03, Geb05\] for details. The $`n`$-braid group $`B_n`$ has the group presentation
$$B_n=\sigma _1,\mathrm{},\sigma _{n1}|\begin{array}{cc}\sigma _i\sigma _j=\sigma _j\sigma _i\hfill & \text{if }|ij|2,\hfill \\ \sigma _i\sigma _j\sigma _i=\sigma _j\sigma _i\sigma _j\hfill & \text{if }|ij|=1.\hfill \end{array},$$
where $`\sigma _i`$ is the isotopy class of the positive half Dehn-twist along the straight line segment connecting the punctures $`i`$ and $`i+1`$. An $`n`$-braid can be regarded as a collection of $`n`$ strands $`l=l_1\mathrm{}l_n`$ in $`[0,1]\times D^2`$ such that $`|l(\{t\}\times D^2)|=n`$ for $`0t1`$ and $`l(\{0,1\}\times D^2)=\{0,1\}\times \{1,\mathrm{},n\}`$.
### 2.1. Positive braid monoid
Let $`B_n^+`$ be the monoid generated by $`\sigma _1,\mathrm{},\sigma _{n1}`$ with the defining relations: $`\sigma _i\sigma _j=\sigma _j\sigma _i`$ for $`|ij|2`$; $`\sigma _i\sigma _j\sigma _i=\sigma _j\sigma _i\sigma _j`$ for $`|ij|=1`$. $`B_n^+`$ is a (left and right) cancellative monoid that embeds in $`B_n`$ under the canonical homomorphism. $`B_n^+`$ is called the *positive braid monoid* and its elements are called *positive braids*.
###### Definition 2.1.
The partial orders $`_L`$ and $`_R`$ on $`B_n^+`$ are defined as follows: for $`P,QB_n^+`$, $`P_LQ`$ if $`Q=PS`$ for some $`SB_n^+`$; $`P_RQ`$ if $`Q=SP`$ for some $`SB_n^+`$.
It is known that the posets $`(B_n^+,_L)`$ and $`(B_n^+,_R)`$ are lattices. Let $`_L`$ and $`_L`$ (respectively, $`_R`$ and $`_R`$) denote the gcd and the lcm with respect to $`_L`$ (respectively, $`_R`$). For positive braids $`P_1`$ and $`P_2`$, the gcd $`P_1_RP_2`$ and the lcm $`P_1_RP_2`$ are characterized by the following properties:
* $`P_1=Q_1(P_1_RP_2)`$ and $`P_2=Q_2(P_1_RP_2)`$ for some $`Q_1,Q_2B_n^+`$ with $`Q_1_RQ_2=1`$;
* $`P_1_RP_2=R_1P_1=R_2P_2`$ for some $`R_1,R_2B_n^+`$ with $`R_1_LR_2=1`$.
The partial orders $`_L`$ and $`_R`$, and thus the lattice structures in $`B_n^+`$ can be extended to $`B_n`$ as follows: for $`\alpha ,\beta B_n`$, $`\alpha _L\beta `$ if $`\beta =\alpha P`$ for some $`PB_n^+`$; $`\alpha _R\beta `$ if $`\beta =P\alpha `$ for some $`PB_n^+`$.
###### Definition 2.2.
The braid $`\mathrm{\Delta }=(\sigma _1\mathrm{}\sigma _{n1})(\sigma _1\mathrm{}\sigma _{n2})\mathrm{}(\sigma _1\sigma _2)\sigma _1`$ is called the *fundamental braid* (or the *Garside element*). Let $`๐=\{AB_n^+:A_L\mathrm{\Delta }\}`$. The elements of $`๐`$ are called *permutation braids* (or *simple elements*).
The fundamental braid $`\mathrm{\Delta }`$ has the following properties: $`A_L\mathrm{\Delta }`$ if and only if $`A_R\mathrm{\Delta }`$ for $`AB_n^+`$; $`\mathrm{\Delta }_LP`$ if and only if $`\mathrm{\Delta }_RP`$ for $`PB_n^+`$; $`\sigma _i_L\mathrm{\Delta }`$ and $`\sigma _i\mathrm{\Delta }=\mathrm{\Delta }\sigma _{ni}`$ for $`i=1,\mathrm{},n1`$. Permutation $`n`$-braids are in one-to-one correspondence with $`n`$-permutations: for an $`n`$-permutation $`\theta `$, the diagram (in $`[0,1]\times `$) of the corresponding braid is obtained by connecting $`(1,i)\{1\}\times `$ to $`(0,\theta (i))\{0\}\times `$ by a straight line for each $`i=1,\mathrm{},n`$ and then making the $`i`$-th strand lie above the $`j`$-th strand whenever $`i<j`$.
For $`PB_n^+`$, let $`\mathrm{s}_L(P)=P_L\mathrm{\Delta }`$ and $`\mathrm{s}_R(P)=P_R\mathrm{\Delta }`$. It is known that for $`P,QB_n^+`$,
$$\mathrm{s}_L(PQ)=\mathrm{s}_L(P\mathrm{s}_L(Q))\text{and}\mathrm{s}_R(PQ)=\mathrm{s}_R(\mathrm{s}_R(P)Q).$$
For $`\alpha B_n`$, there are integers $`uv`$ such that $`\mathrm{\Delta }^u_L\alpha _L\mathrm{\Delta }^v`$. Let $`inf(\alpha )=\mathrm{max}\{u:\mathrm{\Delta }^u_L\alpha \}`$ and $`sup(\alpha )=\mathrm{min}\{v:\alpha _L\mathrm{\Delta }^v\}`$.
###### Definition 2.3.
The expression $`\mathrm{\Delta }^uA_1\mathrm{}A_m`$ is called the *left (respectively, right) normal form* of $`\alpha `$ if $`u=inf(\alpha )`$, $`A_i๐\{1,\mathrm{\Delta }\}`$ and $`\mathrm{s}_L(A_i\mathrm{}A_m)=A_i`$ (respectively, $`\mathrm{s}_R(A_1\mathrm{}A_i)=A_i`$) for $`i=1,\mathrm{},m`$.
###### Definition 2.4.
For $`PB_n^+`$, the *starting set* $`S(P)`$ and the *finishing set* $`F(P)`$ of $`P`$ are defined as
$$S(P)=\{i\sigma _i_LP\}\text{and}F(P)=\{i\sigma _i_RP\}.$$
The following properties are well known \[Thu92, EM94\].
###### Lemma 2.5.
* For a positive braid $`P`$, $`S(\mathrm{s}_L(P))=S(P)`$.
* If $`A`$ is a permutation braid with induced permutation $`\theta `$,
$$S(A)=\{i\theta ^1(i)>\theta ^1(i+1)\}\text{and}F(A)=\{i\theta (i)>\theta (i+1)\}.$$
* For permutation braids $`A`$ and $`B`$, the expression $`AB`$ is in left (respectively, right) normal form if and only if $`F(A)S(B)`$ (respectively, $`F(A)S(B)`$).
By Thurston \[Thu92\], an $`n`$-braid $`\alpha `$ has a unique expression
$$\alpha =P^1Q,$$
where $`P,QB_n^+`$ and $`P_LQ=1`$. We call it the *np*-form of $`\alpha `$. Similarly, we define the *pn*-form of $`\alpha `$ as $`\alpha =PQ^1`$, where $`P,QB_n^+`$ and $`P_RQ=1`$.
Let $`\tau `$ be the inner automorphism of $`B_n`$ defined by $`\tau (\sigma _i)=\sigma _{ni}`$. Then $`\mathrm{\Delta }^1\alpha \mathrm{\Delta }=\tau (\alpha )`$ for $`\alpha B_n`$. The following is known \[Cha95, Lemma 2.3\].
###### Lemma 2.6.
Let $`P,QB_n^+`$. For $`A๐`$, let $`\overline{A}=\mathrm{\Delta }A^1`$.
1. Let $`P=A_mA_{m1}\mathrm{}A_1`$ and $`Q=A_{m+1}A_{m+2}\mathrm{}A_l`$ be in left normal forms. If $`P^1Q`$ is in $`np`$-form, then $`\mathrm{\Delta }^m\tau ^{1m}(\overline{A}_1)\mathrm{}\tau ^1(\overline{A}_{m1})\overline{A}_mA_{m+1}\mathrm{}A_l`$ is the left normal form of $`P^1Q`$.
2. Let $`P=A_1A_2\mathrm{}A_m`$ and $`Q=A_lA_{l1}\mathrm{}A_{m+1}`$ be in right normal forms. If $`PQ^1`$ is in $`pn`$-form, then $`\mathrm{\Delta }^{ml}\tau ^{ml}(A_1)\mathrm{}\tau ^{ml}(A_m)\tau ^{ml+1}(\overline{A}_{m+1})\mathrm{}\tau ^1(\overline{A}_{l1})\overline{A}_l`$ is the right normal form of $`PQ^1`$.
### 2.2. Conjugacy problem in braid groups
Let $`\mathrm{\Delta }^uA_1\mathrm{}A_m`$ be the left normal form of $`\alpha B_n`$. The *cycling* $`๐(\alpha )`$ and the *decycling* $`๐(\alpha )`$ are defined by
$`๐(\alpha )`$ $`=`$ $`\mathrm{\Delta }^uA_2\mathrm{}A_m\tau ^u(A_1);`$
$`๐(\alpha )`$ $`=`$ $`\mathrm{\Delta }^u\tau ^u(A_m)A_1\mathrm{}A_{m1}.`$
Let $`[\alpha ]`$ denote the conjugacy class of $`\alpha `$. Let $`inf_s(\alpha )=\mathrm{max}\{inf(\beta ):\beta [\alpha ]\}`$ and $`sup_s(\alpha )=\mathrm{min}\{sup(\beta ):\beta [\alpha ]\}`$.
###### Definition 2.7.
For $`\alpha B_n`$, the *super summit set* $`[\alpha ]^S`$, the *ultra summit set* $`[\alpha ]^U`$ and the *stable super summit set* $`[\alpha ]^{St}`$ of $`\alpha `$ are defined as follows:
$`[\alpha ]^S`$ $`=`$ $`\{\beta [\alpha ]:inf(\beta )=\underset{s}{inf}(\alpha ),sup(\beta )=\underset{s}{sup}(\alpha )\};`$
$`[\alpha ]^U`$ $`=`$ $`\{\beta [\alpha ]^S:๐^m(\beta )=\beta \text{ for some }m1\};`$
$`[\alpha ]^{St}`$ $`=`$ $`\{\beta [\alpha ]^S:\beta ^m[\alpha ^m]^S\text{for all }m1\}.`$
By definition, $`[\alpha ]^U`$ and $`[\alpha ]^{St}`$ are subsets of $`[\alpha ]^S`$.
###### Theorem 2.8.
Let $`\alpha B_n`$.
1. If $`๐^m(\alpha )=\alpha `$ for some $`m1`$, then $`inf(\alpha )=inf_s(\alpha )`$.
2. If $`๐^m(\alpha )=\alpha `$ for some $`m1`$, then $`sup(\alpha )=sup_s(\alpha )`$.
3. $`๐^m๐^l(\alpha )[\alpha ]^U`$ for some $`m,l0`$.
4. Both $`[\alpha ]^S`$ and $`[\alpha ]^U`$ are finite and nonempty.
5. If $`\beta [\alpha ]^S`$, then $`๐(\beta ),๐(\beta ),\tau (\beta )[\alpha ]^S`$. The same is true for $`[\alpha ]^U`$.
6. If $`\beta [\alpha ]^S`$, then $`๐(๐(\alpha ))=๐(๐(\alpha ))`$, $`\tau (๐(\beta ))=๐(\tau (\beta ))`$ and $`\tau (๐(\beta ))=๐(\tau (\beta ))`$.
7. For $`\beta ,\beta ^{}[\alpha ]^S`$, there is a finite sequence
$$\beta =\beta _0\beta _1\mathrm{}\beta _m=\beta ^{}$$
such that for $`i=0,\mathrm{},m1`$, $`\beta _i[\alpha ]^S`$ and $`\beta _{i+1}=A_i\beta _iA_i^1`$ for some $`A_i๐`$. The same is true for $`[\alpha ]^U`$.
For the results on stable super summit sets, see \[LL06a, LL06b\]. For $`\beta [\alpha ]^S`$, let
$`C^S(\beta )`$ $`=`$ $`\{PB_n^+:P^1\beta P[\beta ]^S\};`$
$`C^U(\beta )`$ $`=`$ $`\{PB_n^+:P^1\beta P[\beta ]^U\}.`$
Both $`C^S(\beta )`$ and $`C^U(\beta )`$ are closed under $`_L`$ by Franco and Gonzรกlez-Meneses \[FG03\] and Gebhardt \[Geb05\], respectively. The closedness under $`_L`$ makes the conjugacy algorithm more efficient.
For a nonempty subset $`๐ฑ`$ of $`B_n^+`$, we call an element $`P๐ฑ`$ the *$`_R`$-minimal element* of $`๐ฑ`$ if $`P_RQ`$ for all $`Q๐ฑ`$. By definition, the $`_R`$-minimal element is unique if it exists. If $`๐ฑ`$ is closed under $`_R`$, then $`๐ฑ`$ has the $`_R`$-minimal element.
The following notions are useful in studying powers \[LL07a, LL06b\]. For $`\alpha B_n`$, let
$$t_{inf}(\alpha )=\underset{m\mathrm{}}{lim}\frac{inf(\alpha ^m)}{m}\text{and}t_{sup}(\alpha )=\underset{m\mathrm{}}{lim}\frac{sup(\alpha ^m)}{m}.$$
The following lists important properties of $`t_{inf}()`$ and $`t_{sup}()`$. See Lemmas 3.2, 3.3, Theorem 3.13 in \[LL07a\], and Corollary 3.5 in \[LL06b\].
###### Theorem 2.9.
Let $`\alpha B_n`$.
1. $`t_{inf}(\gamma \alpha \gamma ^1)=t_{inf}(\alpha )`$ and $`t_{sup}(\gamma \alpha \gamma ^1)=t_{sup}(\alpha )`$ for all $`\gamma B_n`$.
2. $`t_{inf}(\alpha ^m)=mt_{inf}(\alpha )`$ and $`t_{sup}(\alpha ^m)=mt_{sup}(\alpha )`$ for all $`m1`$.
3. $`inf_s(\alpha )t_{inf}(\alpha )<inf_s(\alpha )+1`$ and $`sup_s(\alpha )1<t_{sup}(\alpha )sup_s(\alpha )`$.
4. $`t_{inf}(\alpha )`$ and $`t_{sup}(\alpha )`$ are rational of the form $`p/q`$ for some integers $`p,q`$ with $`1qn(n1)/2`$.
### 2.3. Duality between cycling and decycling
In many aspects, the cycling and the decycling are dual to each other. We define a variant of the cycling as follows so that the duality is more clear. See Lemmas 2.11 and 2.13.
###### Definition 2.10.
For $`\alpha B_n`$, define $`๐_0(\alpha )=\tau ^1(๐(\alpha ))`$.
Since $`\tau ^2(\beta )=\beta `$ and $`\tau (๐(\beta ))=๐(\tau (\beta ))`$ for $`\beta [\alpha ]^S`$, we can replace $`๐`$ with $`๐_0`$ in Theorem 2.8 and in the definition of $`[\alpha ]^U`$. In particular, for an element $`\beta [\alpha ]^S`$, $`\beta `$ belongs to the ultra summit set $`[\alpha ]^U`$ if and only if $`๐_0^m(\beta )=\beta `$ for some $`m1`$.
###### Lemma 2.11.
Let $`\mathrm{\Delta }^uA_1\mathrm{}A_m`$ be the left normal form of $`\alpha B_n`$.
1. The set $`\{PB_n^+:inf(P\alpha )>inf(\alpha )\}`$ is nonempty and closed under $`_R`$. The $`_R`$-minimal element $`A`$ of this set is the permutation braid $`\tau ^u(\mathrm{\Delta }A_1^1)`$ and satisfies $`๐_0(\alpha )=A\alpha A^1`$.
2. The set $`\{PB_n^+:sup(\alpha P^1)<sup(\alpha )\}`$ is nonempty and closed under $`_R`$. The $`_R`$-minimal element $`A`$ of this set is the permutation braid $`A_m`$ and satisfies $`๐(\alpha )=A\alpha A^1`$.
###### Proof.
We prove only (i) since (ii) can be proved similarly. Nonemptiness of $`\{PB_n^+:inf(P\alpha )>inf(\alpha )\}`$ is clear. Note that
* $`(\beta \alpha )_R(\gamma \alpha )=(\beta _R\gamma )\alpha `$ for all $`\alpha ,\beta ,\gamma B_n`$;
* $`inf(\alpha _R\beta )=\mathrm{min}\{inf(\alpha ),inf(\beta )\}`$ for all $`\alpha ,\beta B_n`$.
If $`inf(P\alpha )>inf(\alpha )`$ and $`inf(Q\alpha )>inf(\alpha )`$ for positive braids $`P`$ and $`Q`$, then
$$inf((P_RQ)\alpha )=inf((P\alpha )_R(Q\alpha ))=\mathrm{min}\{inf(P\alpha ),inf(Q\alpha )\}>inf(\alpha ).$$
Therefore, the set $`\{PB_n^+:inf(P\alpha )>inf(\alpha )\}`$ is closed under $`_R`$.
It is easy to see that the $`_R`$-minimal element $`A`$ is $`\tau ^u(\mathrm{\Delta }A_1^1)`$ and, hence,
$`A\alpha A^1`$ $`=`$ $`(\mathrm{\Delta }\tau ^u(A_1^1))(\mathrm{\Delta }^uA_1\mathrm{}A_m)(\tau ^u(A_1)\mathrm{\Delta }^1)`$
$`=`$ $`\mathrm{\Delta }(\mathrm{\Delta }^uA_2\mathrm{}A_m\tau ^u(A_1))\mathrm{\Delta }^1=\mathrm{\Delta }๐(\alpha )\mathrm{\Delta }^1=\tau ^1(๐(\alpha ))`$
$`=`$ $`๐_0(\alpha ).`$
###### Definition 2.12.
For $`\alpha B_n`$, the set
$$[\alpha ]_๐^U=\{\beta [\alpha ]^S:๐^m(\beta )=\beta \text{ for some }m1\}$$
is called the *ultra summit set of $`\alpha `$ with respect to decycling*.
The following lemma is easy to prove, so we omit the proof. It shows that there is a duality between $`๐_0()๐()`$, $`inf()sup()`$ and $`[]^U[]_๐^U`$.
###### Lemma 2.13.
Let $`\alpha B_n`$.
1. $`inf(\alpha )=sup(\alpha ^1)`$ and $`inf_s(\alpha )=sup_s(\alpha ^1)`$.
2. $`๐_0(\alpha )=(๐(\alpha ^1))^1`$.
3. $`\beta [\alpha ]^S`$ if and only if $`\beta ^1[\alpha ^1]^S`$.
4. $`\beta [\alpha ]^U`$ if and only if $`\beta ^1[\alpha ^1]_๐^U`$.
## 3. Braids sending a standard curve to a standard curve
In this section we study the normal form of braids that send a standard curve system to a standard curve system. We collect basic properties of such braids in Lemma 3.5, from which the other results of this section follow easily.
We start by defining some notions. Throughout the paper, we do not distinguish the curves and the isotopy classes of curves.
###### Definition 3.1.
A curve system means a finite collection of disjoint simple closed curves. A simple closed curve in $`D_n`$ is said to be *essential* if it is homotopic neither to a point nor to a puncture nor to the boundary. An essential curve system in $`D_n`$ is said to be *standard* if each component is isotopic to a round circle centered at the real axis as in Figure 1. It is said to be *unnested* if none of its components encloses another component. See Figure 2.
###### Definition 3.2.
The $`n`$-braid group $`B_n`$ acts on the set of curve systems in $`D_n`$. Let $`\alpha ๐`$ denote the left action of $`\alpha B_n`$ on the curve system $`๐`$ in $`D_n`$. An $`n`$-braid $`\alpha `$ is said to be *reducible* if $`\alpha ๐=๐`$ for some essential curve system $`๐`$ in $`D_n`$. Such a curve system $`๐`$ is called a *reduction system* of $`\alpha `$.
The unnested standard curve systems in $`D_n`$ are in one-to-one correspondence with the $`k`$-compositions of $`n`$ for $`2kn1`$. Recall that an ordered $`k`$-tuple $`๐ง=(n_1,\mathrm{},n_k)`$ is a $`k`$-composition of $`n`$ if $`n_i1`$ for each $`i`$ and $`n=n_1+\mathrm{}+n_k`$.
###### Definition 3.3.
For a composition $`๐ง=(n_1,\mathrm{},n_k)`$ of $`n`$, let $`๐_๐ง`$ denote the curve system $`_{n_i2}C_i`$, where $`C_i`$ is the standard curve enclosing $`\{m:_{j=1}^{i1}n_j<m_{j=1}^in_j\}`$. See Figure 2.
The $`k`$-braid group $`B_k`$ acts on the set of $`k`$-compositions of $`n`$ via the induced permutations: for a $`k`$-composition $`๐ง=(n_1,\mathrm{},n_k)`$ and $`\alpha _0B_k`$ with induced permutation $`\theta `$, $`\alpha _0๐ง=(n_{\theta ^1(1)},\mathrm{},n_{\theta ^1(k)})`$.
###### Definition 3.4.
Let $`๐ง=(n_1,\mathrm{},n_k)`$ be a composition of $`n`$.
* Let $`\alpha _0=l_1\mathrm{}l_k`$ be a $`k`$-braid with $`l_i(\{1\}\times D^2)=\{(1,i)\}`$ for each $`i`$. Note that the strands of $`\alpha _0`$ are numbered from bottom to top at its right end. We define $`\alpha _0_๐ง`$ as the $`n`$-braid obtained from $`\alpha _0`$ by taking $`n_i`$ parallel copies of $`l_i`$ for each $`i`$.
* Let $`\alpha _iB_{n_i}`$ for $`i=1,\mathrm{},k`$. We define $`(\alpha _1\mathrm{}\alpha _k)`$ as the $`n`$-braid $`\alpha _1^{}\alpha _2^{}\mathrm{}\alpha _k^{}`$, where $`\alpha _i^{}`$ is the image of $`\alpha _i`$ under the homomorphism $`B_{n_i}B_n`$ defined by $`\sigma _j\sigma _{n_1+\mathrm{}+n_{i1}+j}`$.
We will use the notation $`\alpha =\alpha _0_๐ง(\alpha _1\mathrm{}\alpha _k)`$ throughout the paper. See Figure 3.
###### Lemma 3.5.
Let $`๐ง=(n_1,\mathrm{},n_k)`$ be a composition of $`n`$.
1. The expression $`\alpha =\alpha _0_๐ง(\alpha _1\mathrm{}\alpha _k)`$ is unique, i.e., if $`\alpha _0_๐ง(\alpha _1\mathrm{}\alpha _k)=\beta _0_๐ง(\beta _1\mathrm{}\beta _k)`$, then $`\alpha _i=\beta _i`$ for $`i=0,\mathrm{},k`$.
2. If $`\alpha =\alpha _0_๐ง(\alpha _1\mathrm{}\alpha _k)`$, then $`\alpha ๐_๐ง`$ is standard and, further, $`\alpha ๐_๐ง=๐_{\alpha _0๐ง}`$. Conversely, if $`\alpha ๐_๐ง`$ is standard, then $`\alpha `$ can be expressed as $`\alpha =\alpha _0_๐ง(\alpha _1\mathrm{}\alpha _k)`$.
3. Let $`\alpha =\alpha _0_๐ง(\alpha _1\mathrm{}\alpha _k)`$. If all $`\alpha _i`$โs are positive (respectively, permutation and fundamental) braids, then so is $`\alpha `$.
4. $`\alpha _0_๐ง(\alpha _1\mathrm{}\alpha _k)=(\alpha _{\theta ^1(1)}\mathrm{}\alpha _{\theta ^1(k)})\alpha _0_๐ง`$, where $`\theta `$ is the induced permutation of $`\alpha _0`$.
5. $`\alpha _0\beta _0_๐ง=\alpha _0_{\beta _0๐ง}\beta _0_๐ง`$.
6. $`(\alpha _0_๐ง)^1=\alpha _0^1_{\alpha _0๐ง}`$.
7. $`(\alpha _1\beta _1\mathrm{}\alpha _k\beta _k)=(\alpha _1\mathrm{}\alpha _k)(\beta _1\mathrm{}\beta _k)`$
8. $`(\alpha _1\mathrm{}\alpha _k)^1=(\alpha _1^1\mathrm{}\alpha _k^1)`$.
9. Let $`A_0`$ and $`B_0`$ be permutation $`k`$-braids. $`A_0B_0`$ is in left (respectively, right) normal form if and only if $`A_0_{B_0๐ง}B_0_๐ง`$ is in left (respectively, right) normal form.
10. Let $`P_i`$, $`i=0,\mathrm{},k`$, be positive braids with appropriate braid indices. Let $`A_i=\mathrm{s}_L(P_i)`$ and $`B_i=\mathrm{s}_R(P_i)`$ for $`i=0,\mathrm{},k`$. Then
$`\mathrm{s}_L((P_1\mathrm{}P_k)P_0_๐ง)`$ $`=`$ $`(A_1\mathrm{}A_k)A_0_{(A_0^1P_0)๐ง};`$
$`\mathrm{s}_R(P_0_๐ง(P_1\mathrm{}P_k))`$ $`=`$ $`B_0_๐ง(B_1\mathrm{}B_k).`$
###### Proof.
The statements from (i) to (viii) are easy to prove. Let us prove (ix) and (x).
(ix) Let $`B_0๐ง=(n_1^{},\mathrm{},n_k^{})`$ and $`N_i=n_1^{}+\mathrm{}+n_i^{}`$ for $`i=1,\mathrm{},k`$. Then,
$`F(A_0_{B_0๐ง})`$ $`=`$ $`\{N_i:iF(A_0)\};`$
$`S(B_0_๐ง)`$ $`=`$ $`\{N_i:iS(B_0)\}.`$
Hence, $`F(A_0)S(B_0)`$ if and only if $`F(A_0_{B_0๐ง})S(B_0_๐ง)`$, and $`F(A_0)S(B_0)`$ if and only if $`F(A_0_{B_0๐ง})S(B_0_๐ง)`$.
(x) We prove only the second identity. The first one can be proved in a similar way. It is easy to see that $`\mathrm{s}_R(P_0_๐ง)=B_0_๐ง`$ by (ix) and that $`\mathrm{s}_R(P_1\mathrm{}P_k)=(B_1\mathrm{}B_k)`$. Let $`\theta `$ be the induced permutation of $`B_0`$. Then, by (iv)
$`\mathrm{s}_R(P_0_๐ง(P_1\mathrm{}P_k))=\mathrm{s}_R(\mathrm{s}_R(P_0_๐ง)(P_1\mathrm{}P_k))`$
$`=`$ $`\mathrm{s}_R(B_0_๐ง(P_1\mathrm{}P_k))=\mathrm{s}_R((P_{\theta ^1(1)}\mathrm{}P_{\theta ^1(k)})B_0_๐ง)`$
$`=`$ $`\mathrm{s}_R(\mathrm{s}_R(P_{\theta ^1(1)}\mathrm{}P_{\theta ^1(k)})B_0_๐ง)=\mathrm{s}_R((B_{\theta ^1(1)}\mathrm{}B_{\theta ^1(k)})B_0_๐ง)`$
$`=`$ $`\mathrm{s}_R(B_0_๐ง(B_1\mathrm{}B_k))=B_0_๐ง(B_1\mathrm{}B_k).`$
The last equality holds since $`B_0_๐ง(B_1\mathrm{}B_k)`$ is a permutation braid by (iii). โ
Let $`\mathrm{br}(\alpha )`$ denote the braid index of $`\alpha `$.
###### Lemma 3.6.
Let $`\alpha =\alpha _0_๐ง(\alpha _1\mathrm{}\alpha _k)B_n`$.
1. $`inf(\alpha )=\mathrm{min}\{inf(\alpha _i):i=0,\mathrm{},k,\mathrm{br}(\alpha _i)2\}`$.
2. $`sup(\alpha )=\mathrm{max}\{sup(\alpha _i):i=0,\mathrm{},k,\mathrm{br}(\alpha _i)2\}`$.
3. $`\alpha `$ is a positive (respectively, permutation and fundamental) braid if and only if each $`\alpha _i`$ is a positive (respectively, permutation and fundamental) braid for $`i=0,\mathrm{},k`$.
###### Proof.
(i) Let $`r=\mathrm{min}\{inf(\alpha _i):i=0,\mathrm{},k,\mathrm{br}(\alpha _i)2\}`$. Set $`n_0=k`$. For $`i=0,\mathrm{},k`$, let $`\alpha _i=\mathrm{\Delta }_i^rP_i`$, where $`\mathrm{\Delta }_i`$ is the fundamental braid of $`B_{n_i}`$ and $`P_iB_{n_i}^+`$. Let $`P=P_0_๐ง(P_1\mathrm{}P_k)`$. By Lemma 3.5 (iv), (v) and (vii),
$`\alpha `$ $`=`$ $`\mathrm{\Delta }_0^rP_0_๐ง(\mathrm{\Delta }_1^rP_1\mathrm{}\mathrm{\Delta }_k^rP_k)`$
$`=`$ $`\mathrm{\Delta }_0^r_{P_0๐ง}P_0_๐ง(\mathrm{\Delta }_1^r\mathrm{}\mathrm{\Delta }_k^r)(P_1\mathrm{}P_k)`$
$`=`$ $`\mathrm{\Delta }_0^r_{P_0๐ง}(\mathrm{\Delta }_{\theta ^1(1)}^r\mathrm{}\mathrm{\Delta }_{\theta ^1(k)}^r)P_0_๐ง(P_1\mathrm{}P_k)`$
where $`\theta `$ is the induced permutation of $`P_0`$. Since $`P_0๐ง=(n_{\theta ^1(1)},\mathrm{},n_{\theta ^1(k)})`$, we have $`\mathrm{\Delta }_0^r_{P_0๐ง}(\mathrm{\Delta }_{\theta ^1(1)}^r\mathrm{}\mathrm{\Delta }_{\theta ^1(k)}^r)=\mathrm{\Delta }^r`$, and hence $`\alpha =\mathrm{\Delta }^rP`$. Since $`inf(P_i)=0`$ for some $`P_i`$ with $`\mathrm{br}(P_i)2`$, $`\mathrm{s}_R(P)\mathrm{\Delta }`$ by Lemma 3.5 (x). Therefore $`inf(\alpha )=r`$.
(ii) Since $`sup(\alpha )=inf(\alpha ^1)`$ by Lemma 2.13 (i) and $`\alpha ^1=(\alpha _1^1\mathrm{}\alpha _k^1)\alpha _0^1_{\alpha _0๐ง}`$ by Lemma 3.5 (vi) and (viii), the assertion follows from (i).
(iii) Note that a braid $`\beta `$ is a positive (respectively, permutation and fundamental) braid if and only if $`inf(\beta )0`$ (respectively, $`0inf(\beta )sup(\beta )1`$ and $`inf(\beta )=sup(\beta )=1`$). Therefore, the assertion follows from (i) and (ii) and Lemma 3.5 (iii). โ
###### Lemma 3.7.
Let $`๐`$ be a standard curve system in $`D_n`$ and $`PB_n^+`$ such that $`P๐`$ is standard.
1. If $`P=QA`$ and $`A=\mathrm{s}_R(P)`$, then $`A๐`$ is standard.
2. If $`P=AQ`$ and $`A=\mathrm{s}_L(P)`$, then $`Q๐`$ is standard.
###### Proof.
A curve system is standard if and only if each of its components is standard. Hence, we may assume that the given standard curve system $`๐`$ is unnested. Let $`๐=๐_๐ง`$ for a composition $`๐ง=(n_1,\mathrm{},n_k)`$ of $`n`$.
(i) $`P=P_0_๐ง(P_1\mathrm{}P_k)`$ for some positive braids $`P_i`$, $`i=0,\mathrm{},k`$, by Lemmas 3.5 (ii) and 3.6 (iii). By Lemma 3.5 (x), $`A=\mathrm{s}_R(P)=\mathrm{s}_R(P_0)_๐ง(\mathrm{s}_R(P_1)\mathrm{}\mathrm{s}_R(P_k))`$. By Lemma 3.5 (ii), $`A๐`$ is standard.
(ii) $`P=(P_1\mathrm{}P_k)P_0_๐ง`$ for some positive braids $`P_i`$, $`i=0,\mathrm{},k`$, by Lemmas 3.5 (ii), (iv) and 3.6 (iii). Let $`A_i=\mathrm{s}_L(P_i)`$ for $`i=0,\mathrm{},k`$. Then $`A=\mathrm{s}_L(P)=(A_1\mathrm{}A_k)A_0_{(A_0^1P_0)๐ง}`$ by Lemma 3.5 (x). By Lemma 3.5 (vi) and (viii), $`A^1=A_0^1_{P_0๐ง}(A_1^1\mathrm{}A_k^1)`$. By Lemma 3.5 (ii),
$`Q๐_๐ง`$ $`=`$ $`(A^1P)๐_๐ง=A^1(P๐_๐ง)=A^1๐_{P_0๐ง}`$
$`=`$ $`\left(A_0^1_{P_0๐ง}(A_1^1\mathrm{}A_k^1)\right)๐_{P_0๐ง}=๐_{(A_0^1P_0)๐ง}.`$
Hence $`Q๐`$ is standard. โ
###### Theorem 3.8.
Let $`๐`$ be a standard curve system in $`D_n`$ and $`\mathrm{\Delta }^uA_1\mathrm{}A_m`$ be the (left or right) normal form of $`\alpha B_n`$. If $`\alpha ๐`$ is standard, then so is $`(A_i\mathrm{}A_m)๐`$ for $`i=1,\mathrm{},m`$.
###### Proof.
It is an immediate consequence of Lemma 3.7, because $`(A_1\mathrm{}A_m)๐=\mathrm{\Delta }^u(\alpha ๐)`$ is standard. โ
Roughly speaking, Theorem 3.8 says that if a braid $`\alpha `$ sends a standard curve system to a standard curve system, then so does each permutation braid in the normal form of $`\alpha `$ as in Figure 4.
###### Corollary 3.9.
Let $`\mathrm{\Delta }^uA_1\mathrm{}A_m`$ be the left normal form of an $`n`$-braid $`\alpha `$. If $`\alpha `$ has a standard reduction system $`๐`$, then $`๐_0(\alpha )`$, $`๐(\alpha )`$ and $`\tau (\alpha )`$ have standard reduction systems $`\tau ^u(\mathrm{\Delta }A_1^1)๐`$, $`A_m๐`$ and $`\mathrm{\Delta }^1๐`$, respectively.
###### Proof.
$`A_m๐`$ is standard by Theorem 3.8. By Lemma 2.11,
$$๐(\alpha )(A_m๐)=(A_m\alpha A_m^1)(A_m๐)=A_m(\alpha ๐)=A_m๐.$$
Therefore $`๐(\alpha )`$ has a standard reduction system $`A_m๐`$. In the same way, $`\tau (\alpha )`$ and $`๐_0(\alpha )`$ have standard reduction systems $`\mathrm{\Delta }^1๐`$ and $`\tau ^u(\mathrm{\Delta }A_1^1)๐`$, respectively. โ
###### Corollary 3.10.
Let $`\alpha `$ be a reducible $`n`$-braid with a reduction system $`๐`$. There exists an element $`\beta `$ of the ultra summit set $`[\alpha ]^U`$ which has a standard reduction system. Precisely, there exists a positive braid $`P`$ such that $`\beta =P\alpha P^1`$ belongs to $`[\alpha ]^U`$ and $`P๐`$ is a standard reduction system of $`\beta `$.
###### Proof.
Let $`P_1`$ be a positive $`n`$-braid such that $`P_1๐`$ is standard. Then $`P_1\alpha P_1^1`$ has the standard reduction system $`P_1๐`$. Take $`l,m0`$ such that $`\beta =๐_0^l๐^m(P_1\alpha P_1^1)`$ belongs to $`[\alpha ]^U`$. Lemma 2.11 and Corollary 3.9 say that if $`\gamma B_n`$ has a standard reduction system $`๐^{}`$, then there are permutation braids $`A_1`$ and $`A_2`$ such that $`๐_0(\gamma )=A_1\gamma A_1^1`$ and $`๐(\gamma )=A_2\gamma A_2^1`$ have standard reduction systems $`A_1๐^{}`$ and $`A_2๐^{}`$, respectively. Hence, we can find a positive $`n`$-braid $`P_2`$ such that $`\beta =P_2(P_1\alpha P_1^1)P_2^1`$ and $`P_2(P_1๐)=(P_2P_1)๐`$ is standard. Let $`P_2P_1=P`$. Then, $`\beta =P\alpha P^1`$ and $`\beta `$ has the standard reduction system $`(P_2P_1)๐=P๐`$. โ
###### Corollary 3.11.
Let $`๐`$ be a standard curve system in $`D_n`$, and let $`\alpha ๐`$ be standard for an $`n`$-braid $`\alpha `$.
1. If $`P^1Q`$ is the np-form of $`\alpha `$, then $`Q๐`$ is standard.
2. If $`PQ^1`$ is the pn-form of $`\alpha `$, then $`Q^1๐`$ is standard.
###### Proof.
By Lemma 2.6 and Theorem 3.8, $`Q๐`$ and $`Q^1๐`$ are standard. โ
We remark that Theorem 3.8 and Corollary 3.10 were obtained also by Benardete, Gutiรฉrrez and Nitecki \[BGN95, Theorems 5.7 and 5.8\], and that these two are enough to solve the reducibility problem because there is an efficient algorithm that decides whether a given braid has a standard reduction system or not and finds one if it has \[BGN93\]. However, Corollary 3.10 guarantees only the *existence* of an element (in the ultra summit set of a reducible braid) that has a standard reduction system. To solve the reducibility problem using only Corollary 3.10, we have to compute all the elements in the ultra summit set.
## 4. Standardizers of curve systems
###### Definition 4.1.
For an essential curve system $`๐`$ in $`D_n`$, we define the *standardizer* of $`๐`$ as the set
$$\mathrm{St}(๐)=\{PB_n^+:P๐\text{ is standard}\}.$$
This section is devoted to the study of properties of standardizers. Clearly, $`\mathrm{St}(๐)`$ is nonempty for any essential curve system $`๐`$. Theorem 4.2 shows that standardizers are sublattices of $`B_n^+`$, hence they have unique $`_R`$-minimal elements. The main result of this section is Theorem 4.9 that for any reduction system $`๐`$ of a reducible braid $`\alpha `$, conjugating $`\alpha `$ by the $`_R`$-minimal element of $`\mathrm{St}(๐)`$ preserves the membership of the super summit set, ultra summit set and stable super summit set. Proposition 4.4 and Corollary 4.5 show that the $`_R`$-minimal element of $`\mathrm{St}(๐)`$ does not entangle any standard curve disjoint from $`๐`$. Proposition 4.8 is a characterization of the $`_R`$-minimal element of $`\mathrm{St}(๐)`$ in terms of normal form and lattice operations.
###### Theorem 4.2.
For an essential curve system $`๐`$ in $`D_n`$, its standardizer $`\mathrm{St}(๐)`$ is closed under $`_R`$ and $`_R`$, and hence a sublattice of $`B_n^+`$. Therefore $`\mathrm{St}(๐)`$ contains a unique $`_R`$-minimal element.
###### Proof.
(See Figure 5.)
Let $`P_1,P_2\mathrm{St}(๐)`$. Let $`P_1=Q_1(P_1_RP_2)`$ and $`P_2=Q_2(P_1_RP_2)`$ for $`Q_1,Q_2B_n^+`$ with $`Q_1_RQ_2=1`$. Then $`P_2=Q_2(P_1_RP_2)=Q_2Q_1^1P_1`$, and $`Q_2Q_1^1`$ is in *pn*-form. Since $`P_1๐`$ and $`P_2๐`$ are standard and
$$P_2๐=(Q_2Q_1^1)(P_1๐),$$
$`Q_1^1(P_1๐)=(P_1_RP_2)๐`$ is standard by Corollary 3.11 (ii).
Let $`P_1_RP_2=R_1P_1=R_2P_2`$ for $`R_1,R_2B_n^+`$ with $`R_1_LR_2=1`$. Then $`R_2^1R_1P_1=P_2`$, and $`R_2^1R_1`$ is the *np*-form. Since $`P_1๐`$ and $`P_2๐`$ are standard and
$$P_2๐=(R_2^1R_1)(P_1๐),$$
$`R_1(P_1๐)=(P_1_RP_2)๐`$ is standard by Corollary 3.11 (i). โ
Let $`๐`$, $`๐_1`$ and $`๐_2`$ be essential curve systems such that $`๐=๐_1๐_2`$. Then $`\mathrm{St}(๐)\mathrm{St}(๐_i)`$ for $`i=1,2`$. Let $`P`$, $`P_1`$ and $`P_2`$ be the $`_R`$-minimal elements of $`\mathrm{St}(๐)`$, $`\mathrm{St}(๐_1)`$ and $`\mathrm{St}(๐_2)`$, respectively. By Theorem 4.2, $`P_1_RP`$ and $`P_2_RP`$, hence $`(P_1_RP_2)_RP`$. One may expect that $`P=P_1_RP_2`$. However, the following example shows that it is not true in general.
###### Example 4.3.
Let $`C_1`$ and $`C_2`$ be the curves in $`D_4`$ as in Figure 6. The $`_R`$-minimal elements of $`\mathrm{St}(C_1)`$, $`\mathrm{St}(C_2)`$ and $`\mathrm{St}(C_1C_2)`$ are $`\sigma _1`$, $`\sigma _3`$ and $`\sigma _2\sigma _1\sigma _3`$, respectively. Note that $`\sigma _2\sigma _1\sigma _3`$ is not equal to $`\sigma _1_R\sigma _3=\sigma _1\sigma _3`$.
The following proposition shows that, when an essential curve $`C`$ in $`D_n`$ is standardized by the action of the $`_R`$-minimal element of $`\mathrm{St}(C)`$, any other standard curve disjoint from $`C`$ remains standard.
###### Proposition 4.4.
Let $`C`$ be an essential simple closed curve in $`D_n`$ and let $`P`$ be the $`_R`$-minimal element of $`\mathrm{St}(C)`$. For any standard curve $`C^{}`$ in $`D_n`$ with $`CC^{}=\mathrm{}`$, the curve $`PC^{}`$ is standard.
###### Proof.
Let $`C^{}`$ be a standard curve which is disjoint from $`C`$ and encloses the punctures $`\{r,r+1,\mathrm{},r+s\}`$. Because $`C`$ and $`C^{}`$ are disjoint, $`C`$ is either inside $`C^{}`$ or outside $`C^{}`$ as Figure 7.
Case 1. $`C`$ is inside $`C^{}`$.
There exists a positive braid $`Q`$ written as a positive word on $`\sigma _r,\mathrm{},\sigma _{r+s1}`$ such that $`QC`$ is standard. Since $`Q\mathrm{St}(C)`$ and $`P`$ is the $`_R`$-minimal element of $`\mathrm{St}(C)`$, we have $`P_RQ`$, hence $`Q=RP`$ for some positive braid $`R`$. In particular, $`P`$ is written as a positive word on $`\sigma _r,\mathrm{},\sigma _{r+s1}`$, and hence $`PC^{}=C^{}`$ is standard.
Case 2. $`C`$ is outside $`C^{}`$.
For a braid diagram $`K`$, let $`c(K)`$ denote the number of crossings in $`K`$. Note that if all the crossings in $`K`$ are positive, then $`K`$ represents a positive braid $`Q`$ with $`|Q|=c(K)`$, where $`|Q|`$ denotes the word length of $`Q`$ with respect to $`\sigma _i`$โs.
*Claim. Let $`C`$ and $`C^{}`$ be essential simple closed curves in $`D_n`$ such that $`C^{}`$ is standard and $`C`$ is outside $`C^{}`$. Let $`P`$ be an element (not necessarily the $`_R`$-minimal element) of $`\mathrm{St}(C)`$. Then there is a positive braid $`Q`$ such that $`|Q||P|`$ and both $`QC`$ and $`QC^{}`$ are standard.*
###### Proof of Claim.
See Figure 8 which illustrates this proof with a simple example. Let $`K=l_1\mathrm{}l_n`$ be a braid diagram of $`P`$ in $`[0,1]\times `$ such that the number of crossings in $`K`$ is exactly $`|P|`$. Here we assume that the right end of $`l_i`$ is $`(1,i)`$ for $`i=1,\mathrm{},n`$. Let $`\{r,r+1,\mathrm{},r+s\}`$ be the set of punctures inside $`C^{}`$. Let $`K^{}=l_rl_{r+1}\mathrm{}l_{r+s}`$ and $`K^{\prime \prime }=KK^{}`$. For $`i=r,\mathrm{},r+s`$, let $`e_i`$ be the number of crossings between $`l_i`$ and $`K^{\prime \prime }`$. Let $`e_{i_0}`$ be the minimum of $`\{e_r,e_{r+1},\mathrm{},e_{r+s}\}`$. Then
$$|P|=c(K)=c(K^{})+c(K^{\prime \prime })+(e_r+\mathrm{}+e_{r+s})c(K^{\prime \prime })+(s+1)e_{i_0}.$$
Let $`L`$ be the braid diagram which is the union of $`K^{\prime \prime }`$ and $`(s+1)`$ parallel copies of $`l_{i_0}`$, and let $`Q`$ be the positive braid represented by $`L`$. Since all the crossings in $`L`$ are positive,
$$|Q|=c(L)=c(K^{\prime \prime })+(s+1)e_{i_0}|P|.$$
By the construction of $`Q`$, both the curves $`QC`$ and $`QC^{}`$ are standard. โ
By the above claim, there exists a positive braid $`Q`$ such that $`|Q||P|`$ and both $`QC`$ and $`QC^{}`$ are standard. Because $`P`$ is the $`_R`$-minimal element of $`\mathrm{St}(C)`$ and $`QC`$ is standard, we have $`P_RQ`$. Since $`|Q||P|`$, we obtain $`P=Q`$, hence $`PC^{}`$ is standard. โ
Proposition 4.4 says that if we standardize the components of a curve system $`๐=C_1\mathrm{}C_k`$ one after another by the $`_R`$-minimal element of the standardizers as follows, then the product of the $`_R`$-minimal elements used in this process is exactly the $`_R`$-minimal element of $`\mathrm{St}(๐)`$.
1. Standardize the first component $`C_1`$ of $`๐`$ using the $`_R`$-minimal element $`P_1`$ of $`\mathrm{St}(C_1)`$. Then $`P_1๐=P_1C_1\mathrm{}P_1C_k`$ and $`P_1C_1`$ is standard.
2. Standardize the second component $`P_1C_2`$ of $`P_1๐`$ by the $`_R`$-minimal element $`P_2`$ of $`\mathrm{St}(P_1C_2)`$. Then the first two components $`(P_2P_1)(C_1C_2)`$ of $`(P_2P_1)๐`$ are standard.
3. Continue the above process. Then $`(P_k\mathrm{}P_1)๐`$ is standard. Corollary 4.5 shows that in fact $`P_k\mathrm{}P_1`$ is the $`_R`$-minimal element of $`\mathrm{St}(๐)`$.
###### Corollary 4.5.
Let $`๐,๐_1,\mathrm{},๐_k`$ be essential curve systems in $`D_n`$ such that $`๐=๐_1\mathrm{}๐_k`$. Let $`P`$ be the $`_R`$-minimal element of $`\mathrm{St}(๐)`$.
1. If $`P_i`$ is the $`_R`$-minimal element of $`\mathrm{St}((P_{i1}\mathrm{}P_1)๐_i)`$, then $`P=P_kP_{k1}\mathrm{}P_1`$.
2. For any standard curve $`C^{}`$ disjoint from $`๐`$, the curve $`PC^{}`$ is standard.
###### Proof.
We prove the corollary only for the case when each curve system $`๐_i`$ has only one component. The general case can be proved easily from this. Suppose that each curve system $`๐_i`$ has only one component.
*Claim.* The following hold for each $`i=0,1,\mathrm{},k`$.
* $`P_iP_{i1}\mathrm{}P_1_RP`$.
* The curve $`(P_iP_{i1}\mathrm{}P_1)๐_j`$ is standard for $`j=1,\mathrm{},i`$.
* For any standard curve $`C^{}`$ disjoint from $`๐`$, the curve $`(P_iP_{i1}\mathrm{}P_1)C^{}`$ is standard.
###### Proof of Claim.
The statement is obvious for $`i=0`$ since $`P_i\mathrm{}P_1`$ is the identity. Using induction on $`i`$, assume that the statement is true for some $`i`$ with $`0i<k`$. Since $`P_i\mathrm{}P_1_RP`$,
$$P=Q(P_i\mathrm{}P_1)$$
for some $`QB_n^+`$. Since $`Q((P_i\mathrm{}P_1)๐_{i+1})=P๐_{i+1}`$ is standard and $`P_{i+1}`$ is the $`_R`$-minimal element of $`\mathrm{St}((P_i\mathrm{}P_1)๐_{i+1})`$, we have $`P_{i+1}_RQ`$, hence
$$P_{i+1}P_i\mathrm{}P_1_RQ(P_i\mathrm{}P_1)=P.$$
By the induction hypothesis, $`(P_i\mathrm{}P_1)C^{}`$ and $`(P_i\mathrm{}P_1)๐_j`$ are standard curves disjoint from $`(P_i\mathrm{}P_1)๐_{i+1}`$ for $`j=1,\mathrm{},i`$. Since $`P_{i+1}`$ is the $`_R`$-minimal element of $`\mathrm{St}((P_i\mathrm{}P_1)๐_{i+1})`$, $`(P_{i+1}P_i\mathrm{}P_1)C^{}`$ and $`(P_{i+1}P_i\mathrm{}P_1)๐_j`$ for $`j=1,\mathrm{},i`$ are standard by Proposition 4.4. By definition of $`P_{i+1}`$, $`(P_{i+1}P_i\mathrm{}P_1)๐_{i+1}`$ is standard. โ
By (b) of the above claim, $`(P_kP_{k1}\mathrm{}P_1)๐`$ is standard. Since $`P`$ is the $`_R`$ minimal element of $`\mathrm{St}(๐)`$, $`P_R(P_kP_{k1}\mathrm{}P_1)`$. By (a) of the claim, $`(P_kP_{k1}\mathrm{}P_1)_RP`$, hence $`P=P_kP_{k1}\mathrm{}P_1`$. By (c) of the claim, $`PC^{}`$ is standard for any standard curve $`C^{}`$ disjoint from $`๐`$. โ
In the rest of this section, we use the following definition.
###### Definition 4.6.
For a composition $`๐ง=(n_1,\mathrm{},n_k)`$ of $`n`$, we define the symbol $`\delta _๐ง`$ and non-negative integers $`N_0,N_1,\mathrm{},N_k`$ as follows:
* $`\delta _๐ง=\mathrm{\Delta }_1\mathrm{}\mathrm{\Delta }_k`$, where $`\mathrm{\Delta }_i`$ is the fundamental braid of $`B_{n_i}`$ for $`i=1,\mathrm{},k`$;
* $`N_0=0`$ and $`N_i=n_1+n_2+\mathrm{}+n_i`$ for $`i=1,\mathrm{},k`$.
Then, for a composition $`๐ง=(n_1,\mathrm{},n_k)`$ of $`n`$ and $`\sigma _iB_k`$, the following hold.
* If $`A_L\delta _๐ง`$, then $`A๐_๐ง=๐_๐ง`$.
* $`S(\delta _๐ง)=F(\delta _๐ง)=\{1,\mathrm{},n1\}\{N_1,\mathrm{},N_{k1}\}`$.
* $`\sigma _i๐ง=\sigma _i^1๐ง=(n_1,\mathrm{},n_{i1},n_{i+1},n_i,n_{i+2},\mathrm{},n_k)`$.
* $`\delta _๐ง\sigma _i_{\sigma _i๐ง}=\sigma _i_{\sigma _i๐ง}\delta _{\sigma _i๐ง}`$. See Figure 9.
###### Lemma 4.7.
Let $`๐ง=(n_1,\mathrm{},n_k)`$ be a composition of $`n`$.
1. Let $`A`$ be a permutation $`n`$-braid with induced permutation $`\theta `$. Then $`\delta _๐งA`$ is a permutation braid if and only if $`\theta ^1`$ is order-preserving on the set $`\{N_{i1}+1,\mathrm{},N_i\}`$ for each $`i=1,\mathrm{},k`$, that is,
$$\theta ^1(N_{i1}+1)<\theta ^1(N_{i1}+2)<\mathrm{}<\theta ^1(N_i).$$
2. For a positive $`n`$-braid $`P`$, the starting set $`S(\delta _๐งP)`$ is strictly greater than the starting set $`S(\delta _๐ง)`$ if and only if $`\sigma _i_{\sigma _i๐ง}_LP`$ for some $`i\{1,\mathrm{},k1\}`$.
###### Proof.
(i) It is an easy consequence of the fact that a positive braid $`P`$ is a permutation braid if and only if any two of its strands cross at most once \[Thu92, Lemma9.1.10\] or \[EM94, Lemma 2.3\]. See Figure 10.
(ii) See Figure 11.
Suppose $`\sigma _i_{\sigma _i๐ง}_LP`$ for some $`i\{1,\mathrm{},k1\}`$. Then $`N_iS(\delta _๐งP)`$, hence $`S(\delta _๐งP)`$ is strictly greater than $`S(\delta _๐ง)`$. Conversely, suppose that $`S(\delta _๐งP)`$ is strictly greater than $`S(\delta _๐ง)`$. Let $`A`$ be the permutation $`n`$-braid such that $`\mathrm{s}_L(\delta _๐งP)=\delta _๐งA`$, that is, $`\delta _๐งA`$ is the first permutation braid in the left normal form of $`\delta _๐งP`$. Then $`N_iS(\delta _๐งA)`$ for some $`i\{1,\mathrm{},k1\}`$. Let $`\omega `$ and $`\theta `$ be the induced permutations of $`\delta _๐ง`$ and $`A`$ respectively. Then
$$\omega ^1(N_i)=N_{i1}+1\text{and}\omega ^1(N_i+1)=N_{i+1}.$$
Since $`N_iS(\delta _๐งA)`$, we have $`(\omega \theta )^1(N_i)>(\omega \theta )^1(N_i+1)`$ and, hence,
(1)
$$\theta ^1(N_{i1}+1)>\theta ^1(N_{i+1}).$$
Because $`\theta ^1`$ is order-preserving on each of the sets $`\{N_{i1}+1,N_{i1}+2,\mathrm{},N_i\}`$ and $`\{N_i+1,N_i+2,\mathrm{},N_{i+1}\}`$, we have the following:
(2) $`\theta ^1(N_{i1}+1)<\mathrm{}<\theta ^1(N_i1)<\theta ^1(N_i);`$
(3) $`\theta ^1(N_i+1)<\theta ^1(N_i+2)<\mathrm{}<\theta ^1(N_{i+1}).`$
From (1), (2) and (3), we obtain $`\sigma _i_{\sigma _i๐ง}_LA_LP`$. โ
The following proposition characterizes the minimal element of the standardizer $`\mathrm{St}(๐)`$ of a curve system $`๐`$.
###### Proposition 4.8.
Let $`๐`$ be an unnested curve system in $`D_n`$. Let $`P`$ be a positive braid such that $`P๐`$ is standard and, hence, $`P๐=๐_๐ง`$ for some composition $`๐ง=(n_1,\mathrm{},n_k)`$ of $`n`$. Then the following conditions are equivalent.
1. $`P`$ is the $`_R`$-minimal element of the standardizer $`\mathrm{St}(๐)`$.
2. $`P_L\delta _๐ง=1`$ and $`S(\delta _๐งP)=S(\delta _๐ง)`$.
3. $`P^1(\delta _๐งP)`$ is in *np*-form.
4. $`P^1(\delta _๐ง^lP)`$ is in *np*-form for some $`l1`$.
5. $`P^1(\delta _๐ง^lP)`$ is in *np*-form for all $`l1`$.
###### Proof.
We prove the equivalence by showing that (i) $``$ (ii) $``$ (v) $``$ (iii) $``$ (iv) $``$ (ii). The implications (v) $``$ (iii) and (iii) $``$ (iv) are obvious.
(i) $``$ (ii) Let $`A=P_L\delta _๐ง`$ and let $`P=AQ`$ for some positive braid $`Q`$. Since $`A_L\delta _๐ง`$, $`A๐_๐ง=๐_๐ง`$, and hence
$$Q๐=A^1(P๐)=A^1๐_๐ง=๐_๐ง.$$
Therefore $`Q\mathrm{St}(๐)`$. By the $`_R`$-minimality of $`P`$, we have $`P=Q`$ and, hence, $`P_L\delta _๐ง=A=1`$.
Assume that $`S(\delta _๐งP)`$ is strictly greater than $`S(\delta _๐ง)`$. Then, by Lemma 4.7 (ii), $`P=\sigma _i_{\sigma _i๐ง}Q`$ for some $`i\{1,\mathrm{},k1\}`$ and some positive braid $`Q`$. Since
$$Q๐=(\sigma _i_{\sigma _i๐ง})^1(P๐)=\sigma _i^1_๐ง๐_๐ง,$$
$`Q๐`$ is standard. This contradicts the $`_R`$-minimality of $`P`$. Consequently, $`S(\delta _๐งP)=S(\delta _๐ง)`$.
(ii) $``$ (i) Let $`Q`$ be the $`_R`$-minimal element of $`\mathrm{St}(๐)`$. Let $`Q๐=๐_๐ง^{}`$ for some composition $`๐ง^{}`$ of $`n`$. Since $`P๐`$ is standard, $`P=RQ`$ for some positive braid $`R`$. Since
$$R๐_๐ง^{}=R(Q๐)=P๐=๐_๐ง,$$
the positive braid $`R`$ sends the standard curve system $`๐_๐ง^{}`$ to the standard curve system $`๐_๐ง`$. Therefore, by Lemmas 3.5 (ii) and 3.6 (iii), $`R=R_0_๐ง^{}(R_1\mathrm{}R_k)`$ for some positive braids $`R_i`$ with appropriate braid indices, and $`R_0๐ง^{}=๐ง`$.
If $`(R_1\mathrm{}R_k)1`$, then $`P_L\delta _๐ง1`$. This contradicts the hypothesis. Therefore $`(R_1\mathrm{}R_k)=1`$.
If $`R_01`$, then $`R_0=\sigma _iR_0^{}`$ for some $`i\{1,\mathrm{},k1\}`$ and a positive $`k`$-braid $`R_0^{}`$. Since $`R_0^{}๐ง^{}=\sigma _i^1(R_0๐ง^{})=\sigma _i^1๐ง=\sigma _i๐ง`$,
$$R_0_๐ง^{}=\sigma _i_{R_0^{}๐ง^{}}R_0^{}_๐ง^{}=\sigma _i_{\sigma _i๐ง}R_0^{}_๐ง^{}.$$
Since $`\sigma _i_{\sigma _i๐ง}_LR_0_๐ง^{}_LP`$, $`S(\delta _๐งP)`$ is strictly greater than $`S(\delta _๐ง)`$ by Lemmas 3.5 (ii). This contradicts the hypothesis $`S(\delta _๐งP)=S(\delta _๐ง)`$. Therefore $`R=1`$ and, hence, $`P`$ is the $`_R`$-minimal element of $`\mathrm{St}(๐)`$.
(ii) $``$ (v) We first claim that $`S(\delta _๐ง^lP)=S(\delta _๐ง)`$ for all $`l1`$. Let $`\delta _๐งA=\mathrm{s}_L(\delta _๐งP)`$. Since $`S(\delta _๐งA)=S(\delta _๐งP)`$ by Lemma 2.5 (i) and $`S(\delta _๐งP)=S(\delta _๐ง)`$ by the hypothesis,
$$S(\delta _๐งA)=S(\delta _๐งP)=S(\delta _๐ง)=F(\delta _๐ง).$$
In particular, $`F(\delta _๐ง)S(\delta _๐งA)`$, and hence $`\delta _๐ง(\delta _๐งA)`$ is in left normal form by Lemma 2.5 (iii). Since $`F(\delta _๐ง)=S(\delta _๐ง)`$, $`\underset{l1}{\underset{}{\delta _๐ง\mathrm{}\delta _๐ง}}(\delta _๐งA)`$ is the left normal form of $`\delta _๐ง^lA`$ for all $`l1`$, and hence $`S(\delta _๐ง^lP)=S(\delta _๐ง^lA)=S(\delta _๐ง)`$.
Now we have $`S(\delta _๐ง^lP)=S(\delta _๐ง)`$ for all $`l1`$. By the hypothesis $`P_L\delta _๐ง=1`$,
$$S(P)S(\delta _๐ง^lP)=S(P)S(\delta _๐ง)=\mathrm{}\text{for all }l1.$$
Consequently, $`P_L\delta _๐ง^lP=1`$ and $`P^1(\delta _๐ง^lP)`$ is in *np*-form for all $`l1`$.
(iv) $``$ (ii) Let $`P^1(\delta _๐ง^lP)`$ is in *np*-form for some $`l1`$, that is, $`P_L(\delta _๐ง^lP)=1`$. Since $`P_L\delta _๐ง_LP_L(\delta _๐ง^lP)`$, we have $`P_L\delta _๐ง=1`$.
Assume that $`S(\delta _๐งP)`$ is strictly greater than $`S(\delta _๐ง)`$. By Lemma 4.7 (ii), we have
(4)
$$P=\sigma _i_{\sigma _i๐ง}Q$$
for some $`i\{1,\mathrm{},k1\}`$ and some positive braid $`Q`$. Since $`\delta _๐ง\sigma _i_{\sigma _i๐ง}=\sigma _i_{\sigma _i๐ง}\delta _{\sigma _i๐ง}`$,
(5)
$$\delta _๐ง^lP=\delta _๐ง^l\sigma _i_{\sigma _i๐ง}Q=\sigma _i_{\sigma _i๐ง}\delta _{\sigma _i๐ง}^lQ.$$
By (4) and (5), we obtain $`\sigma _i_{\sigma _i๐ง}_LP_L(\delta _๐ง^lP)`$, which contracts the hypothesis that $`P^1(\delta _๐ง^lP)`$ is in *np*-from. As a result, $`S(\delta _๐งP)=S(\delta _๐ง)`$. โ
Now we are ready to show that standardizing a reduction system $`๐`$ of a braid by the $`_R`$-minimal element of $`\mathrm{St}(๐)`$ preserves the membership of the super summit set, ultra summit set and stable super summit set. The anonymous referee of this journal pointed out that our initial proof of the following theorem contains a mistake. The proof is corrected as suggested by the referee.
###### Theorem 4.9.
Let $`\alpha `$ be a reducible $`n`$-braid with a reduction system $`๐`$. Let $`P`$ be the $`_R`$-minimal element of $`\mathrm{St}(๐)`$. Then the following hold.
1. $`inf(\alpha )inf(P\alpha P^1)sup(P\alpha P^1)sup(\alpha )`$.
2. If $`\alpha [\alpha ]^S`$, then $`P\alpha P^1[\alpha ]^S`$.
3. If $`\alpha [\alpha ]^U`$, then $`P\alpha P^1[\alpha ]^U`$.
4. If $`\alpha [\alpha ]^{St}`$, then $`P\alpha P^1[\alpha ]^{St}`$.
###### Proof.
First, suppose that $`๐`$ is an unnested curve system. Let $`P๐=๐_๐ง`$ for a composition $`๐ง=(n_1,\mathrm{},n_k)`$ of $`n`$. Let $`u=sup(P)`$. Define $`\overline{P}=\mathrm{\Delta }^uP^1`$ and $`Q=\overline{P}\delta _๐ง^2P`$. By Proposition 4.8, $`P^1(\delta _๐ง^2P)`$ is in *np*-form, hence, by Lemma 2.6 (i),
$$\overline{P}=Q_L\mathrm{\Delta }^{sup(\overline{P})}.$$
Since $`(P\alpha P^1)๐_๐ง=๐_๐ง`$, $`P\alpha P^1=\beta _0_๐ง(\beta _1\mathrm{}\beta _k)`$ for some $`\beta _i`$โs with appropriate braid indices, and $`\beta _0๐ง=๐ง`$. Thus $`P\alpha P^1`$ commutes with $`\delta _๐ง^2`$, and it follows that $`\alpha `$ commutes with $`P^1\delta _๐ง^2P`$. Therefore $`Q\alpha Q^1=\left(\mathrm{\Delta }^uP^1\delta _๐ง^2P\right)\alpha \left(\mathrm{\Delta }^uP^1\delta _๐ง^2P\right)^1=\tau ^u(\alpha )`$. That is,
(6)
$$Q^1\tau ^u(\alpha )Q=\alpha .$$
Consider the following sets:
$`C(\alpha )`$ $`=`$ $`\{RB_n^+:inf(\alpha )inf(R^1\alpha R)sup(R^1\alpha R)sup(\alpha )\};`$
$`C^S(\alpha )`$ $`=`$ $`\{RB_n^+:R^1\alpha R[\alpha ]^S\};`$
$`C^U(\alpha )`$ $`=`$ $`\{RB_n^+:R^1\alpha R[\alpha ]^U\};`$
$`C^{St}(\alpha )`$ $`=`$ $`\{RB_n^+:R^1\alpha R[\alpha ]^{St}\}.`$
By Franco and Gonzรกlez-Meneses \[FG03\], Gebhardt \[Geb05\] and Lee and Lee \[LL06a\], all the sets $`C(\alpha )`$, $`C^S(\alpha )`$, $`C^U(\alpha )`$ and $`C^{St}(\alpha )`$ are closed under $`_L`$.
Suppose $`\alpha [\alpha ]^S`$. Since $`\tau ^m(\alpha )[\alpha ]^S`$ for all $`m`$, $`\mathrm{\Delta }^{sup(\overline{P})}C^S(\tau ^u(\alpha ))`$. Since $`QC^S(\tau ^u(\alpha ))`$ by (6), we have $`\overline{P}=Q_L\mathrm{\Delta }^{sup(\overline{P})}C^S(\tau ^u(\alpha ))`$. That is,
$$P\alpha P^1=\overline{P}^1\mathrm{\Delta }^u\alpha \mathrm{\Delta }^u\overline{P}=\overline{P}^1\tau ^u(\alpha )\overline{P}[\tau ^u(\alpha )]^S=[\alpha ]^S.$$
Hence (ii) is proved. The other statements can be proved similarly.
Now we consider general case. For a reduction system $`๐`$ of $`\alpha `$, we decompose $`๐`$ into $`๐_1\mathrm{}๐_l`$, where $`๐_i`$โs are inductively defined as the outermost component of $`๐(๐_1\mathrm{}๐_{i1})`$. By the construction, $`๐_1,\mathrm{},๐_l`$ are unnested reduction systems of $`\alpha `$. For $`i=1,\mathrm{},l`$, define positive braids $`P_i`$ and conjugates $`\alpha _i`$ of $`\alpha `$ inductively as follows. Let $`P_0=1`$ and $`\alpha _0=\alpha `$.
* $`P_i`$ is the $`_R`$-minimal element of $`\mathrm{St}((P_{i1}\mathrm{}P_1)๐_i)`$;
* $`\alpha _i=P_i\alpha _{i1}P_i^1=(P_i\mathrm{}P_1)\alpha (P_i\mathrm{}P_1)^1`$.
Note that each $`\alpha _i`$ is a reducible braid with a reduction system $`(P_i\mathrm{}P_1)๐_{i+1}`$ and that $`P=P_l\mathrm{}P_1`$ by Corollary 4.5 (i).
Suppose $`\alpha [\alpha ]^S`$. By the previous discussion on the unnested case, $`P_{i+1}\alpha _iP_{i+1}^1[\alpha ]^S`$ for $`i=0,\mathrm{},l1`$, hence $`P\alpha P^1[\alpha ]^S`$. Therefore (ii) is proved. The other statements can be proved similarly. โ
## 5. Outermost components of non-periodic reducible braids
In this section we define the outermost component $`\alpha _{\mathrm{ext}}`$ of a non-periodic reducible braid $`\alpha `$ using the $`_R`$-minimal element of the standardizer of the canonical reduction system of $`\alpha `$, and study its properties.
Recall the canonical reduction system of mapping classes. For a reduction system $`๐D_n`$ of an $`n`$-braid $`\alpha `$, let $`D_๐`$ be the closure of $`D_nN(๐)`$ in $`D_n`$, where $`N(๐)`$ is a regular neighborhood of $`๐`$. The restriction of $`\alpha `$ induces an automorphism on $`D_๐`$ that is well-defined up to isotopy. Due to Birman, Lubotzky and McCarthy \[BLM83\] and Ivanov \[Iva92\], for any $`n`$-braid $`\alpha `$, there is a unique *canonical reduction system* $`(\alpha )`$ with the following properties.
1. $`(\alpha ^m)=(\alpha )`$ for all $`m0`$.
2. $`(\beta \alpha \beta ^1)=\beta (\alpha )`$ for all $`\beta B_n`$.
3. The restriction of $`\alpha `$ to each component of $`D_{(\alpha )}`$ is either periodic or pseudo-Anosov. A reduction system with this property is said to be *adequate*.
4. If $`๐`$ is an adequate reduction system of $`\alpha `$, then $`(\alpha )๐`$.
By the properties of canonical reduction systems, a braid $`\alpha `$ is non-periodic reducible if and only if $`(\alpha )\mathrm{}`$. Let $`_{\mathrm{ext}}(\alpha )`$ denote the collection of the outermost components of $`(\alpha )`$. Then, $`_{\mathrm{ext}}(\alpha )`$ is an unnested curve system satisfying the properties (i) and (ii). We remark that, while the canonical reduction systems are defined for the mapping classes of surfaces with genus, we have to restrict ourselves to the mapping classes of punctured disks in order to define the outermost component $`_{\mathrm{ext}}(\alpha )`$.
###### Lemma 5.1.
Let $`\alpha ,\beta B_n`$ with $`(\alpha )\mathrm{}`$. If $`\alpha \beta =\beta \alpha `$, then $`(\alpha )`$ and $`_{\mathrm{ext}}(\alpha )`$ are reduction systems of $`\beta `$.
###### Proof.
Since $`(\alpha )=(\beta \alpha \beta ^1)=\beta (\alpha )`$ and $`_{\mathrm{ext}}(\alpha )=_{\mathrm{ext}}(\beta \alpha \beta ^1)=\beta _{\mathrm{ext}}(\alpha )`$, both $`(\alpha )`$ and $`_{\mathrm{ext}}(\alpha )`$ are reduction systems of $`\beta `$. โ
###### Definition 5.2.
Let $`\alpha B_n`$ with $`(\alpha )\mathrm{}`$. Let $`P`$ be the $`_R`$-minimal element of $`\mathrm{St}(_{\mathrm{ext}}(\alpha ))`$ and $`\beta =P\alpha P^1`$. Since $`_{\mathrm{ext}}(\beta )`$ is unnested and standard, $`_{\mathrm{ext}}(\beta )=๐_๐ง`$ for a composition $`๐ง=(n_1,\mathrm{},n_k)`$ of $`n`$, and $`\beta `$ has the unique expression $`\beta =\beta _0_๐ง(\beta _1\mathrm{}\beta _k)`$ by Lemma 3.5 (ii). We define the *outermost component* $`\alpha _{\mathrm{ext}}`$ of $`\alpha `$ by $`\alpha _{\mathrm{ext}}=\beta _0`$.
In other words, $`\alpha _{\mathrm{ext}}`$ is the restriction of $`\alpha `$ to the outermost component of $`D_n_{\mathrm{ext}}(\alpha )`$. This element is a priori defined up to conjugacy, but the use of the $`_R`$-minimal element $`P`$ determines the particular element $`\beta _0`$ to be chosen in the conjugacy class.
###### Lemma 5.3.
Let $`\alpha `$ be an $`n`$-braid with $`(\alpha )\mathrm{}`$.
* If $`\beta `$ is conjugate to $`\alpha `$, then $`\beta _{\mathrm{ext}}`$ is conjugate to $`\alpha _{\mathrm{ext}}`$.
* $`(\alpha ^m)_{\mathrm{ext}}=(\alpha _{\mathrm{ext}})^m`$ for all $`m0`$.
* $`inf(\alpha )inf(\alpha _{\mathrm{ext}})sup(\alpha _{\mathrm{ext}})sup(\alpha )`$.
* $`inf_s(\alpha )inf_s(\alpha _{\mathrm{ext}})sup_s(\alpha _{\mathrm{ext}})sup_s(\alpha )`$.
* $`t_{inf}(\alpha )t_{inf}(\alpha _{\mathrm{ext}})t_{sup}(\alpha _{\mathrm{ext}})t_{sup}(\alpha )`$.
###### Proof.
(i) is obvious. (ii) follows from $`(\alpha ^m)=(\alpha )`$. (iii) follows from Lemma 3.6 and Theorem 4.9.
(iv) Choose any $`\beta [\alpha ]^S`$. By (iii), we have
$$\underset{s}{inf}(\alpha )=inf(\beta )inf(\beta _{\mathrm{ext}})sup(\beta _{\mathrm{ext}})sup(\beta )=\underset{s}{sup}(\alpha ).$$
Since $`\alpha _{\mathrm{ext}}`$ and $`\beta _{\mathrm{ext}}`$ are conjugate by (i),
$$inf(\beta _{\mathrm{ext}})\underset{s}{inf}(\alpha _{\mathrm{ext}})\underset{s}{sup}(\alpha _{\mathrm{ext}})sup(\beta _{\mathrm{ext}}).$$
Combining the above two, we obtain $`inf_s(\alpha )inf_s(\alpha _{\mathrm{ext}})sup_s(\alpha _{\mathrm{ext}})sup_s(\alpha )`$.
(v) By (ii) and (iii), for all $`m1`$,
$$inf(\alpha ^m)inf((\alpha ^m)_{\mathrm{ext}})=inf((\alpha _{\mathrm{ext}})^m)sup((\alpha _{\mathrm{ext}})^m)=sup((\alpha ^m)_{\mathrm{ext}})sup(\alpha ^m).$$
Therefore,
$$\frac{inf(\alpha ^m)}{m}\frac{inf((\alpha _{\mathrm{ext}})^m)}{m}\frac{sup((\alpha _{\mathrm{ext}})^m)}{m}\frac{sup(\alpha ^m)}{m}.$$
By taking $`m\mathrm{}`$, we obtain the desired inequalities for $`t_{inf}()`$ and $`t_{sup}()`$. โ
###### Lemma 5.4.
Let $`\alpha B_n`$ with $`_{\mathrm{ext}}(\alpha )`$ standard. Then $`_{\mathrm{ext}}(\tau (\alpha ))`$, $`_{\mathrm{ext}}(๐_0(\alpha ))`$ and $`_{\mathrm{ext}}(๐(\alpha ))`$ are standard. Moreover,
1. $`\tau (\alpha )_{\mathrm{ext}}=\tau (\alpha _{\mathrm{ext}});`$
2. $`๐_0(\alpha )_{\mathrm{ext}}=\{\begin{array}{cc}\alpha _{\mathrm{ext}}\hfill & \text{if }inf(\alpha _{\mathrm{ext}})>inf(\alpha );\hfill \\ ๐_0(\alpha _{\mathrm{ext}})\hfill & \text{if }inf(\alpha _{\mathrm{ext}})=inf(\alpha );\hfill \end{array}`$
3. $`๐(\alpha )_{\mathrm{ext}}=\{\begin{array}{cc}\alpha _{\mathrm{ext}}\hfill & \text{if }sup(\alpha _{\mathrm{ext}})<sup(\alpha );\hfill \\ ๐(\alpha _{\mathrm{ext}})\hfill & \text{if }sup(\alpha _{\mathrm{ext}})=sup(\alpha )\text{.}\hfill \end{array}`$
###### Proof.
$`_{\mathrm{ext}}(\tau (\alpha ))=_{\mathrm{ext}}(\mathrm{\Delta }^1\alpha \mathrm{\Delta })=\mathrm{\Delta }^1_{\mathrm{ext}}(\alpha )`$ is obviously standard. $`_{\mathrm{ext}}(๐_0(\alpha ))`$ and $`_{\mathrm{ext}}(๐(\alpha ))`$ are standard by Corollary 3.9. Let $`_{\mathrm{ext}}(\alpha )=๐_๐ง`$ for a composition $`๐ง=(n_1,\mathrm{},n_k)`$ of $`n`$ and $`\alpha =\alpha _0_๐ง(\alpha _1\mathrm{}\alpha _k)`$. Let $`\mathrm{\Delta }_i`$ be the fundamental braid of $`B_{n_i}`$ for $`i=1,\mathrm{},k`$ and $`\mathrm{\Delta }_0`$ be the fundamental braid of $`B_k`$. Note that $`\alpha _0๐ง=๐ง`$ and
$$\mathrm{\Delta }=(\mathrm{\Delta }_1\mathrm{}\mathrm{\Delta }_k)\mathrm{\Delta }_0_{\mathrm{\Delta }_0^1๐ง}=\mathrm{\Delta }_0_{\mathrm{\Delta }_0^1๐ง}(\mathrm{\Delta }_k\mathrm{}\mathrm{\Delta }_1).$$
Therefore,
$`\tau (\alpha )`$ $`=`$ $`\mathrm{\Delta }^1\alpha \mathrm{\Delta }`$
$`=`$ $`\mathrm{\Delta }_0^1_๐ง(\mathrm{\Delta }_1^1\mathrm{}\mathrm{\Delta }_k^1)\alpha _0_๐ง(\alpha _1\mathrm{}\alpha _k)\mathrm{\Delta }_0_{\mathrm{\Delta }_0^1๐ง}(\mathrm{\Delta }_k\mathrm{}\mathrm{\Delta }_1)`$
$`=`$ $`\mathrm{\Delta }_0^1\alpha _0\mathrm{\Delta }_0_{\mathrm{\Delta }_0^1๐ง}(\mathrm{\Delta }_k^1\alpha _k\mathrm{\Delta }_k\mathrm{}\mathrm{\Delta }_1^1\alpha _1\mathrm{\Delta }_1)`$
$`=`$ $`\tau (\alpha _0)_{\mathrm{\Delta }_0^1๐ง}(\tau (\alpha _k)\mathrm{}\tau (\alpha _1)).`$
Since $`_{\mathrm{ext}}(\tau (\alpha ))=\mathrm{\Delta }^1๐_๐ง=๐_{\mathrm{\Delta }_0^1๐ง}`$, $`\tau (\alpha )_{\mathrm{ext}}=\tau (\alpha _0)=\tau (\alpha _{\mathrm{ext}})`$.
Let $`\alpha =\mathrm{\Delta }^uA_1\mathrm{}A_l`$ be the left normal form of $`\alpha `$. Since $`\alpha ๐_๐ง=๐_๐ง`$ is standard, $`A_l๐_๐ง`$ is standard by Theorem 3.8. By Lemmas 3.5 (ii) and 3.6 (iii), $`A_l`$ is expressed as $`A_l=A_{l,0}_๐ง(A_{l,1}\mathrm{}A_{l,k})`$, where $`A_{l,i}`$โs are permutation $`n_i`$-braids. Let $`\theta _1`$ and $`\theta _2`$ be the induced permutations of $`\alpha _0A_{l,0}^1`$ and $`A_{l,0}^1`$ respectively. Then
$`๐(\alpha )`$ $`=`$ $`A_l\alpha A_l^1`$
$`=`$ $`A_{l,0}_๐ง(A_{l,1}\mathrm{}A_{l,k})\alpha _0_๐ง(\alpha _1\mathrm{}\alpha _k)(A_{l,1}^1\mathrm{}A_{l,k}^1)A_{l,0}^1_{A_{l,0}๐ง}`$
$`=`$ $`A_{l,0}\alpha _0A_{l,0}^1_{A_{l,0}๐ง}(A_{l,\theta _1(1)}\alpha _{\theta _2(1)}A_{l,\theta _2(1)}^1\mathrm{}A_{l,\theta _1(k)}\alpha _{\theta _2(k)}A_{l,\theta _2(k)}^1).`$
Recall Lemma 3.6 (ii) that $`sup(\alpha _{\mathrm{ext}})sup(\alpha )`$. If $`sup(\alpha _{\mathrm{ext}})<sup(\alpha )`$, then $`A_{l,0}=1`$ and, hence, $`๐(\alpha )_{\mathrm{ext}}=\alpha _{\mathrm{ext}}`$. If $`sup(\alpha _{\mathrm{ext}})=sup(\alpha )`$, then $`A_{l,0}1`$ and, hence, $`๐(\alpha )_{\mathrm{ext}}=A_{l,0}\alpha _0A_{l,0}^1=๐(\alpha _{\mathrm{ext}})`$.
For $`๐_0(\alpha )`$, use the identity $`๐_0(\alpha )=๐(\alpha ^1)^1`$. โ
Recall Lemma 5.3 that $`inf(\alpha )inf(\alpha _{\mathrm{ext}})`$ and $`inf_s(\alpha )inf_s(\alpha _{\mathrm{ext}})`$ for any $`\alpha B_n`$ with $`(\alpha )\mathrm{}`$.
###### Lemma 5.5.
Let $`\alpha `$ be an $`n`$-braid with $`(\alpha )\mathrm{}`$. Let $`\beta `$ be an element of $`[\alpha ]^U`$ with $`_{\mathrm{ext}}(\beta )`$ standard.
* Let $`inf_s(\alpha _{\mathrm{ext}})>inf_s(\alpha )`$. Then, $`inf(\beta _{\mathrm{ext}})>inf(\beta )`$.
* Let $`inf_s(\alpha _{\mathrm{ext}})=inf_s(\alpha )`$. Then, $`inf(\beta _{\mathrm{ext}})=inf(\beta )`$, and $`๐_0^m(\beta _{\mathrm{ext}})=\beta _{\mathrm{ext}}`$ for some $`m1`$.
###### Proof.
By Lemma 5.3 (i), $`\beta _{\mathrm{ext}}`$ and $`\alpha _{\mathrm{ext}}`$ are conjugate, hence $`inf(\beta _{\mathrm{ext}})inf_s(\alpha _{\mathrm{ext}})`$.
We first prove the following claim.
*Claim*. Let $`inf(\beta _{\mathrm{ext}})=inf(\beta )`$. Then, $`๐_0^m(\beta _{\mathrm{ext}})=\beta _{\mathrm{ext}}`$ for some $`m1`$, and $`inf_s(\alpha _{\mathrm{ext}})=inf(\beta _{\mathrm{ext}})=inf(\beta )=inf_s(\alpha )`$.
###### Proof of Claim.
By Lemma 5.4 (ii), the sequence $`\{inf(๐_0^i(\beta )_{\mathrm{ext}})\}_{i=0}^{\mathrm{}}`$ is non-decreasing. Since $`\beta [\alpha ]^U`$, $`๐_0^m(\beta )=\beta `$ for some $`m1`$. Therefore,
$$inf(๐_0^i(\beta )_{\mathrm{ext}})=inf(\beta _{\mathrm{ext}})\text{for all }i0.$$
Since $`๐_0^i(\beta )[\alpha ]^U`$ for all $`i0`$, we have $`inf(๐_0^i(\beta ))=inf_s(\alpha )=inf(\beta )`$ for all $`i0`$. Hence
$$inf(๐_0^i(\beta )_{\mathrm{ext}})=inf(\beta _{\mathrm{ext}})=inf(\beta )=inf(๐_0^i(\beta ))\text{for all }i0.$$
By Lemma 5.4 (ii),
$$๐_0^i(\beta )_{\mathrm{ext}}=๐_0^i(\beta _{\mathrm{ext}})\text{for all }i0.$$
Since $`๐_0^m(\beta )=\beta `$, we obtain $`๐_0^m(\beta _{\mathrm{ext}})=๐_0^m(\beta )_{\mathrm{ext}}=\beta _{\mathrm{ext}}`$.
By Theorem 2.8 (i), $`inf(\beta _{\mathrm{ext}})=inf_s(\beta _{\mathrm{ext}})=inf_s(\alpha _{\mathrm{ext}})`$. Therefore, $`inf_s(\alpha _{\mathrm{ext}})=inf(\beta _{\mathrm{ext}})=inf(\beta )=inf_s(\alpha )`$. โ
(i) Assume $`inf(\beta _{\mathrm{ext}})=inf(\beta )`$. Then $`inf_s(\alpha _{\mathrm{ext}})=inf(\beta _{\mathrm{ext}})=inf(\beta )=inf_s(\alpha )`$ by the above claim. This contradicts the hypothesis that $`inf_s(\alpha _{\mathrm{ext}})>inf_s(\alpha )`$, hence $`inf(\beta _{\mathrm{ext}})>inf(\beta )`$.
(ii) Since $`inf(\beta )inf(\beta _{\mathrm{ext}})inf_s(\alpha _{\mathrm{ext}})`$,
$$inf(\beta )inf(\beta _{\mathrm{ext}})\underset{s}{inf}(\alpha _{\mathrm{ext}})=\underset{s}{inf}(\alpha )=inf(\beta ).$$
Therefore $`inf(\beta _{\mathrm{ext}})=inf(\beta )`$. By the claim, $`๐_0^m(\beta _{\mathrm{ext}})=\beta _{\mathrm{ext}}`$ for some $`m1`$. โ
The following proposition show that the property $`inf_s(\alpha _{\mathrm{ext}})>inf_s(\alpha )`$ is preserved by taking powers.
###### Proposition 5.6.
If $`inf_s(\alpha _{\mathrm{ext}})>inf_s(\alpha )`$, then $`inf_s((\alpha ^m)_{\mathrm{ext}})>inf_s(\alpha ^m)`$ for all $`m1`$.
###### Proof.
By Theorem 6.1 in \[Lee07\], for any $`\beta B_n`$ and any $`m1`$,
$$\underset{s}{inf}(\beta )\frac{inf_s(\beta ^m)}{m}<\underset{s}{inf}(\beta )+1.$$
By Lemma 5.3 (ii), $`(\alpha ^m)_{\mathrm{ext}}=(\alpha _{\mathrm{ext}})^m`$ for all $`m0`$. Since $`inf_s(\alpha _{\mathrm{ext}})>inf_s(\alpha )`$,
$$\frac{inf_s(\alpha ^m)}{m}<\underset{s}{inf}(\alpha )+1\underset{s}{inf}(\alpha _{\mathrm{ext}})\frac{inf_s((\alpha _{\mathrm{ext}})^m)}{m}=\frac{inf_s((\alpha ^m)_{\mathrm{ext}})}{m}$$
for all $`m1`$. Therefore $`inf_s((\alpha ^m)_{\mathrm{ext}})>inf_s(\alpha ^m)`$ for all $`m1`$. โ
## 6. Split braids
An $`n`$-braid $`\alpha `$ is called a *split braid* if it is conjugate to an element in the subgroup of $`B_n`$ generated by $`\sigma _1,\mathrm{},\sigma _{i1},\sigma _{i+1},\mathrm{},\sigma _{n1}`$ for some $`1in1`$ \[Hum91\]. In our terminology, $`\alpha B_n`$ is a split braid if it is conjugate to a braid $`\beta `$ of the form $`\beta =1_๐ง(\beta _1\beta _2)`$, where $`๐ง=(i,ni)`$ for some $`1in1`$, and $`\beta _1B_i`$ and $`\beta _2B_{ni}`$.
The following lemma is easy to show, but we include a proof for completeness.
###### Lemma 6.1.
Let $`\alpha `$ be an $`n`$-braid.
* $`\alpha `$ is a split braid if and only if either $`\alpha `$ is the identity or $`\alpha `$ is non-periodic and reducible with $`\alpha _{\mathrm{ext}}=1`$.
* Let $`\alpha =1_๐ง(\alpha _1\mathrm{}\alpha _k)`$ for a composition $`๐ง=(n_1,\mathrm{},n_k)`$ of $`n`$ such that $`_{\mathrm{ext}}(\alpha )\mathrm{}`$. Then $`_{\mathrm{ext}}(\alpha )=๐_๐ง`$ if and only if $`\alpha _i`$ is non-split for each $`1ik`$.
###### Proof.
For unnested curve systems $`๐`$ and $`๐^{}`$ in $`D_n`$, let us write โ$`๐^{}๐`$โ if each component of $`๐^{}`$ is enclosed by (possibly parallel to) a component of $`๐`$, and write โ$`๐^{}๐`$โ if $`๐^{}๐`$ and $`๐^{}๐`$. Then $``$ is a partial order over the set of unnested curve systems in $`D_n`$. For compositions $`๐ง=(n_1,\mathrm{},n_k)`$ and $`๐ง^{}`$ of $`n`$, $`๐_๐ง^{}๐_๐ง`$ if and only if $`๐ง^{}`$ is a refinement of $`๐ง`$, that is, for each $`1ik`$, there exists a composition $`(n_{i,1}^{},\mathrm{},n_{i,r_i}^{})`$ of $`n_i`$ such that $`๐ง^{}=(n_{1,1}^{},\mathrm{},n_{1,r_1}^{},\mathrm{},n_{k,1}^{},\mathrm{},n_{k,r_k}^{})`$.
*Claim*. Let $`\beta ๐_๐ง=๐_๐ง`$ for a composition $`๐ง`$ of $`n`$, then $`\beta `$ is written as $`\beta =\beta _0_๐ง(\beta _1\mathrm{}\beta _k)`$. If $`\beta _0`$ is periodic or pseudo-Anosov, then $`_{\mathrm{ext}}(\beta )๐_๐ง`$, and there exists $`PB_n^+`$ such that both $`P_{\mathrm{ext}}(\beta )`$ and $`P๐_๐ง`$ are standard.
###### Proof of Claim.
Because $`\beta _0`$ is periodic or pseudo-Anosov, we can make an adequate reduction system of $`\beta `$ from $`๐_๐ง`$ by adding some curves each of which is enclosed by a component of $`๐_๐ง`$. Because any adequate reduction system of $`\beta `$ contains $`_{\mathrm{ext}}(\beta )`$ as a subset, we have $`_{\mathrm{ext}}(\beta )๐_๐ง`$. Let $`P`$ be the $`_R`$-minimal element of $`\mathrm{St}(_{\mathrm{ext}}(\beta ))`$. Then $`P_{\mathrm{ext}}(\beta )`$ is standard by the construction. Apply Corollary 4.5 to $`๐_๐ง\backslash _{\mathrm{ext}}(\beta )`$. Then $`P(๐_๐ง\backslash _{\mathrm{ext}}(\beta ))`$ and hence $`P๐_๐ง`$ are standard. โ
(i) It is obvious that if $`\alpha `$ is the identity or $`\alpha `$ is non-periodic and reducible with $`\alpha _{\mathrm{ext}}=1`$ then $`\alpha `$ is a split braid. Conversely, suppose that $`\alpha `$ is a split braid. Taking a conjugate of $`\alpha `$ if necessary, we may assume that
$$\alpha =1_๐ง(\alpha _1\alpha _2),$$
where $`๐ง=(\mathrm{},n\mathrm{})`$ for some $`1\mathrm{}n1`$, and $`\alpha _1B_{\mathrm{}}`$ and $`\alpha _2B_n\mathrm{}`$.
First, assume that $`_{\mathrm{ext}}(\alpha )=\mathrm{}`$, that is, $`\alpha `$ is periodic or pseudo-Anosov. Since split braids are a special type of reducible braids and since pseudo-Anosov braids cannot be reducible \[FLP79\], $`\alpha `$ is periodic. Therefore $`\alpha ^p=\mathrm{\Delta }^{2m}`$ for some $`p0`$ and $`m`$, hence
$$1_๐ง(\alpha _1^p\alpha _2^p)=\alpha ^p=\mathrm{\Delta }^{2m}=\mathrm{\Delta }_0^{2m}_๐ง(\mathrm{\Delta }_1^{2m}\mathrm{\Delta }_2^{2m}),$$
where $`\mathrm{\Delta }_0`$, $`\mathrm{\Delta }_1`$ and $`\mathrm{\Delta }_2`$ are the fundamental braids of $`B_2`$, $`B_{\mathrm{}}`$ and $`B_n\mathrm{}`$, respectively. By Lemma 3.5 (i), we have $`\mathrm{\Delta }_0^{2m}=1`$, hence $`m=0`$, and it follows that $`\alpha ^p=1`$. Because braid groups are torsion-free \[Deh98\], $`\alpha `$ is the identity.
Now, assume that $`_{\mathrm{ext}}(\alpha )\mathrm{}`$, that is, $`\alpha `$ is non-periodic and reducible. For a curve system $`๐`$ in a punctured disk $`D_m`$, let $`\mathrm{Out}(D_m๐)`$ denote the outermost component of $`D_m๐`$. By the above claim, we have $`_{\mathrm{ext}}(\alpha )๐_๐ง`$ and hence $`๐_๐ง\mathrm{Out}(D_n_{\mathrm{ext}}(\alpha ))`$, and we may assume that $`_{\mathrm{ext}}(\alpha )`$ is standard. Let $`\alpha _{\mathrm{ext}}`$ be a $`k`$-braid. Because $`_{\mathrm{ext}}(\alpha )`$ is standard, $`\mathrm{Out}(D_n_{\mathrm{ext}}(\alpha ))`$ is canonically diffeomorphic to $`D_k`$. Let $`๐^{}`$ be the image of $`๐_๐ง`$ under this diffeomorphism. Then $`๐^{}`$ is a reduction system of $`\alpha _{\mathrm{ext}}`$ such that the restriction of $`\alpha _{\mathrm{ext}}`$ to $`\mathrm{Out}(D_k๐^{})`$ is the same as the restriction of $`\alpha `$ to $`\mathrm{Out}(D_n๐_๐ง)`$ which is the identity. This means that $`\alpha _{\mathrm{ext}}`$ is a split braid. Because $`\alpha _{\mathrm{ext}}`$ is either periodic or pseudo-Anosov, the discussion in the above paragraph shows that $`\alpha _{\mathrm{ext}}`$ is the identity.
(ii) Assume that $`\alpha _{\mathrm{}}`$ is a split braid for some $`1\mathrm{}k`$, hence $`\alpha _{\mathrm{}}`$ is conjugate to $`1_๐ง_{\mathrm{}}(\alpha _{\mathrm{}}^{}\alpha _{\mathrm{}}^{\prime \prime })`$, where $`๐ง_{\mathrm{}}=(n_{\mathrm{}}^{},n_{\mathrm{}}^{\prime \prime })`$ is a composition of $`n_{\mathrm{}}`$, and $`\alpha _{\mathrm{}}^{}B_n_{\mathrm{}}^{}`$ and $`\alpha _{\mathrm{}}^{\prime \prime }B_{n_{\mathrm{}}^{\prime \prime }}`$. By taking a conjugate of $`\alpha `$ if necessary, we may assume that
$$\alpha =1_๐ง^{}(\alpha _1\mathrm{}\alpha _\mathrm{}1\alpha _{\mathrm{}}^{}\alpha _{\mathrm{}}^{\prime \prime }\alpha _{\mathrm{}+1}\mathrm{}\alpha _k),$$
where $`๐ง^{}=(n_1,\mathrm{},n_\mathrm{}1,n_{\mathrm{}}^{},n_{\mathrm{}}^{\prime \prime },n_{\mathrm{}+1},\mathrm{},n_k)`$. Note that $`๐_๐ง^{}๐_๐ง`$. By the claim, $`_{\mathrm{ext}}(\alpha )๐_๐ง^{}๐_๐ง`$, hence $`๐_๐ง_{\mathrm{ext}}(\alpha )`$.
Conversely, assume that $`_{\mathrm{ext}}(\alpha )๐_๐ง`$. By the claim, we may assume that $`_{\mathrm{ext}}(\alpha )๐_๐ง`$ and $`_{\mathrm{ext}}(\alpha )`$ is standard. Let $`_{\mathrm{ext}}(\alpha )=๐_๐ง^{}`$ for a composition $`๐ง^{}`$ of $`n`$. Then $`๐ง^{}`$ is a refinement of $`๐ง`$, hence, for each $`i`$, there exists a composition $`(n_{i,1}^{},\mathrm{},n_{i,r_i}^{})`$ of $`n_i`$ such that $`๐ง^{}=(n_{1,1}^{},\mathrm{},n_{1,r_1}^{},\mathrm{},n_{k,1}^{},\mathrm{},n_{k,r_k}^{})`$. Because $`\alpha _{\mathrm{ext}}`$ is the identity by (i), $`\alpha `$ is written as
$$\alpha =1_๐ง^{}(\underset{r_1}{\underset{}{\alpha _{1,1}\mathrm{}\alpha _{1,r_1}}}\mathrm{}\underset{r_k}{\underset{}{\alpha _{k,1}\mathrm{}\alpha _{k,r_k}}}).$$
Since $`๐_๐ง^{}=_{\mathrm{ext}}(\alpha )๐_๐ง`$ by the assumption, we have $`r_{\mathrm{}}2`$ for some $`1\mathrm{}k`$. Comparing the above expression with $`\alpha =1_๐ง(\alpha _1\mathrm{}\alpha _k)`$, we have $`\alpha _{\mathrm{}}=1_๐ง_{\mathrm{}}^{}(\alpha _{\mathrm{},1}\mathrm{}\alpha _{\mathrm{},r_{\mathrm{}}})`$ where $`๐ง_{\mathrm{}}^{}=(n_{\mathrm{},1}^{},\mathrm{},n_{\mathrm{},r_{\mathrm{}}}^{})`$. Since $`r_{\mathrm{}}2`$, $`\alpha _{\mathrm{}}`$ is a split braid. โ
For $`\alpha B_n`$, let $`|\alpha |`$ denote the minimal word length of $`\alpha `$ with respect to $`\{\sigma _1^{\pm 1},\mathrm{},\sigma _{n1}^{\pm 1}\}`$. Then $`|\alpha |`$ is the minimum number of crossings in the braid diagram of $`\alpha `$.
###### Proposition 6.2.
If $`\alpha 1`$ is a split braid and $`|\alpha |`$ is minimal in the conjugacy class of $`\alpha `$, then $`_{\mathrm{ext}}(\alpha )`$ is standard.
###### Proof.
There exists a braid $`\beta `$ in the conjugacy class of $`\alpha `$ such that $`_{\mathrm{ext}}(\beta )`$ is standard. Therefore $`_{\mathrm{ext}}(\beta )=๐_๐ง`$ for some composition $`๐ง=(n_1,\mathrm{},n_k)`$ of $`n`$. Then by Lemma 6.1
$$\beta =1_๐ง(\beta _1\mathrm{}\beta _k)$$
for some non-split $`n_i`$-braids $`\beta _i`$. We may choose $`\beta `$ so that $`|\beta _i|`$ is minimal in the conjugacy class of $`\beta _i`$ for each $`i\{1,\mathrm{},k\}`$.
Since $`\alpha `$ and $`\beta `$ are conjugate, $`\alpha =\gamma \beta \gamma ^1`$ for some $`\gamma B_n`$. Let $`\theta `$ be the induced permutation of $`\gamma `$. For $`i=1,\mathrm{},k`$, let $`S_i=\{j:n_1+\mathrm{}+n_{i1}<jn_1+\mathrm{}+n_i\}`$ and $`T_i=\{\theta (j):jS_i\}`$. Let $`\gamma _i`$ be the result of forgetting the $`j`$-th strand from $`\gamma `$ for all $`jS_i`$. (The strands of a braid are numbered from bottom to top at its right end.) See Figure 12. Let $`\alpha _i`$ be the result of forgetting the $`j`$-th strand from $`\alpha `$ for all $`jT_i`$. Then $`\alpha _i=\gamma _i\beta _i\gamma _i^1`$ for all $`i=1,\mathrm{},k`$.
Let $`K`$ be a braid diagram of $`\alpha `$ such that the number of crossings in $`K`$ is exactly $`|\alpha |`$. For $`i=1,\mathrm{},k`$, let $`K_i`$ be the result of deleting the $`j`$-th strand from $`K`$ for all $`jT_i`$. Then $`K_i`$ is a braid diagram of $`\alpha _i`$ for all $`i`$. Let $`c(K)`$ and $`c(K_i)`$ denote the numbers of crossings in $`K`$ and $`K_i`$, respectively. Then $`|\alpha |=c(K)`$, $`|\alpha _i|c(K_i)`$ for each $`i`$ and $`_{i=1}^kc(K_i)c(K)`$.
Since $`|\alpha |`$ is minimal in the conjugacy class, $`|\alpha ||\beta |`$. Since $`|\beta _i|`$ is minimal in the conjugacy class, $`|\beta _i||\alpha _i|`$ for all $`i=1,\mathrm{},k`$. Hence
$$c(K)=|\alpha ||\beta |=\underset{i=1}{\overset{k}{}}|\beta _i|\underset{i=1}{\overset{k}{}}|\alpha _i|\underset{i=1}{\overset{k}{}}c(K_i)c(K).$$
Therefore $`c(K)=_{i=1}^kc(K_i)`$ and it follows that there is no crossing between the strands in $`K_i`$ and those in $`K_j`$ whenever $`ij`$.
Now we claim that each $`T_l`$ is a set of consecutive integers. On the contrary, assume that there exists $`jT_m`$ for some $`ml`$ such that $`i_1<j<i_2`$ for some $`i_1,i_2T_l`$. Let $`K_{l,1}`$ be the result of deleting all $`i`$-th strands from $`K_l`$ with $`i>j`$ and let $`K_{l,2}=K_lK_{l,1}`$. See Figure 13. Because there is no crossing between the strands in $`K_l`$ and those in $`K_m`$, there is no crossing between $`K_{l,1}`$ and $`K_{l,2}`$. Therefore $`K_l`$ is splitted into $`K_{l,1}`$ and $`K_{l,2}`$. This contradicts that $`\alpha _l`$ is non-split. Hence, each $`T_l`$ is a set of consecutive integers.
Let $`T_{i_1},T_{i_2},\mathrm{},T_{i_k}`$ be the rearrangement of $`T_j`$โs such that the elements of the sets are increasing, and let $`๐ง^{}=(n_{i_1},\mathrm{},n_{i_k})`$. Then $`\alpha =1_๐ง^{}(\alpha _{i_1}\mathrm{}\alpha _{i_k})`$ and $`_{\mathrm{ext}}(\alpha )=๐_๐ง^{}`$. Therefore $`_{\mathrm{ext}}(\alpha )`$ is standard. โ
###### Corollary 6.3.
If $`P1`$ is a positive split braid, then $`_{\mathrm{ext}}(P)`$ is standard.
###### Proof.
If $`P`$ is a positive braid, then $`|P|`$ is minimal in the conjugacy class of $`P`$. โ
## 7. Ultra summit sets of reducible braids
In this section, we establish Theorem 7.4, the main result of this paper. Roughly speaking, it says that if the outermost component $`\alpha _{\mathrm{ext}}`$ is simpler than the whole braid $`\alpha `$ from a Garside-theoretic point of view, then it is easy to find a reduction system of $`\alpha `$.
###### Definition 7.1.
Let $`\alpha B_n`$, $`\beta [\alpha ]^U`$ and $`m=\mathrm{min}\{l1:๐_0^l(\beta )=\beta \}`$. For $`i=0,\mathrm{},m1`$, let $`A_i`$ be the $`_R`$-minimal element of $`\{PB_n^+:inf(P๐_0^i(\beta ))>inf(๐_0^i(\beta ))\}`$. The product $`A_{m1}A_{m2}\mathrm{}A_0`$ is called the *cycling commutator* of $`\beta `$ and denoted $`T_\beta `$.
By definition, the cycling commutator $`T_\beta `$ is a positive braid. By Lemma 2.11 (i),
$`T_\beta \beta T_\beta ^1`$ $`=`$ $`A_{m1}\mathrm{}A_2A_1A_0\beta A_0^1A_1^1A_2^1\mathrm{}A_{m1}^1`$
$`=`$ $`A_{m1}\mathrm{}A_2A_1๐_0(\beta )A_1^1A_2^1\mathrm{}A_{m1}^1`$
$`=`$ $`A_{m1}\mathrm{}A_2๐_0^2(\beta )A_2^1\mathrm{}A_{m1}^1`$
$`=`$ $`\mathrm{}=๐_0^m(\beta )=\beta .`$
###### Lemma 7.2.
Let $`\alpha B_n`$ and $`\beta [\alpha ]^U`$. Then the cycling commutator $`T_\beta `$ is a non-identity positive braid with $`T_\beta \beta =\beta T_\beta `$.
The following proposition is the key to Theorem 7.4. We prove it in ยง8.
###### Proposition 7.3.
Let $`\alpha `$ be a non-periodic reducible $`n`$-braid with $`inf_s(\alpha _{\mathrm{ext}})>inf_s(\alpha )`$. For any element $`\beta `$ of $`[\alpha ]^U`$, the cycling commutator $`T_\beta `$ is a split braid.
Recall from Lemma 5.3 that $`inf_s(\alpha )inf_s(\alpha _{\mathrm{ext}})sup_s(\alpha _{\mathrm{ext}})sup_s(\alpha )`$ and $`t_{inf}(\alpha )t_{inf}(\alpha _{\mathrm{ext}})t_{sup}(\alpha _{\mathrm{ext}})t_{sup}(\alpha )`$ for any non-periodic reducible braid $`\alpha `$.
###### Theorem 7.4.
Let $`\alpha `$ be a non-periodic reducible $`n`$-braid.
1. If $`inf_s(\alpha _{\mathrm{ext}})>inf_s(\alpha )`$, then each element of $`[\alpha ]^U`$ has a standard reduction system.
2. If $`sup_s(\alpha _{\mathrm{ext}})<sup_s(\alpha )`$, then each element of $`[\alpha ]_๐^U`$ has a standard reduction system.
3. If $`\alpha `$ is a split braid, then each element of $`[\alpha ]^U[\alpha ]_๐^U`$ has a standard reduction system.
4. If $`\alpha _{\mathrm{ext}}`$ is periodic, then there exists $`1q<n`$ such that each element of $`[\alpha ^q]^U[\alpha ^q]_๐^U`$ has a standard reduction system.
5. If $`t_{inf}(\alpha _{\mathrm{ext}})>t_{inf}(\alpha )`$, then there exists $`1q<n(n1)/2`$ such that each element of $`[\alpha ^q]^U`$ has a standard reduction system.
6. If $`t_{sup}(\alpha _{\mathrm{ext}})<t_{sup}(\alpha )`$, then there exists $`1q<n(n1)/2`$ such that each element of $`[\alpha ^q]_๐^U`$ has a standard reduction system.
###### Proof.
(i) Let $`\beta `$ be an element of $`[\alpha ]^U`$. By Proposition 7.3, the cycling commutator $`T_\beta `$ is a non-identity positive split braid. By Corollary 6.3, $`_{\mathrm{ext}}(T_\beta )`$ is standard. Since $`\beta `$ commutes with $`T_\beta `$ by Lemma 7.2, $`_{\mathrm{ext}}(T_\beta )`$ is a standard reduction system of $`\beta `$ by Lemma 5.1.
(ii) Because $`inf_s((\alpha ^1)_{\mathrm{ext}})=inf_s((\alpha _{\mathrm{ext}})^1)=sup_s(\alpha _{\mathrm{ext}})`$ and $`inf_s(\alpha ^1)=sup_s(\alpha )`$, we have $`inf_s((\alpha ^1)_{\mathrm{ext}})>inf_s(\alpha ^1)`$. By (i), each element of $`[\alpha ^1]^U`$ has a standard reduction system. Because $`[\alpha ]_๐^U=\{\beta ^1:\beta [\alpha ^1]^U\}`$, we are done.
(iii) Let $`\beta [\alpha ]^U`$. If $`inf_s(\alpha )<inf_s(\alpha _{\mathrm{ext}})`$, then $`\beta `$ has a standard reduction system by (i). If $`inf_s(\alpha )=inf_s(\alpha _{\mathrm{ext}})`$, then $`inf(\beta )=inf_s(\alpha )=inf_s(\alpha _{\mathrm{ext}})=0`$ and, hence, $`\beta `$ is positive. Since $`\beta `$ is split, $`_{\mathrm{ext}}(\beta )`$ is standard by Corollary 6.3.
Since $`\alpha `$ is a split braid, so is $`\alpha ^1`$. Thus, every element of $`[\alpha ^1]^U`$ and, hence, $`[\alpha ]_๐^U`$ has a standard reduction system.
(iv) Let $`k`$ be the braid index of $`\alpha _{\mathrm{ext}}`$. Because $`\alpha _{\mathrm{ext}}`$ is periodic, there exist $`1qk`$ and $`l`$ such that
$$(\alpha _{\mathrm{ext}})^q=\mathrm{\Delta }_0^{2l},$$
where $`\mathrm{\Delta }_0`$ is the fundamental braid of $`B_k`$. Then $`\mathrm{\Delta }^{2l}\alpha ^q1`$ is a split braid. By (iii), every element of $`[\mathrm{\Delta }^{2l}\alpha ^q]^U[\mathrm{\Delta }^{2l}\alpha ^q]_๐^U`$ has a standard reduction system. Since
$$[\alpha ^q]^U=\{\mathrm{\Delta }^{2l}\beta :\beta [\mathrm{\Delta }^{2l}\alpha ^q]^U\}\text{and}[\alpha ^q]_๐^U=\{\mathrm{\Delta }^{2l}\beta :\beta [\mathrm{\Delta }^{2l}\alpha ^q]_๐^U\},$$
each element of $`[\alpha ^q]^U[\alpha ^q]_๐^U`$ has a standard reduction system.
(v) Recall from Theorem 2.9 that, for any $`\gamma B_n`$,
* $`t_{inf}(\gamma )`$ is rational with denominator less than or equal to $`|\mathrm{\Delta }|=n(n1)/2`$;
* $`inf_s(\gamma )t_{inf}(\gamma )<inf_s(\gamma )+1`$;
* $`t_{inf}(\gamma ^m)=mt_{inf}(\gamma )`$ for all integers $`m1`$.
Let $`k`$ be the braid index of $`\alpha _{\mathrm{ext}}`$. Then $`t_{inf}(\alpha _{\mathrm{ext}})=p/q`$ for some integers $`p`$, $`q`$ with $`1qk(k1)/2`$. Since $`t_{inf}((\alpha _{\mathrm{ext}})^q)=qt_{inf}(\alpha _{\mathrm{ext}})`$ is an integer, we have $`inf_s((\alpha _{\mathrm{ext}})^q)=qt_{inf}(\alpha _{\mathrm{ext}})`$. Therefore,
$$\underset{s}{inf}((\alpha ^q)_{\mathrm{ext}})=\underset{s}{inf}((\alpha _{\mathrm{ext}})^q)=qt_{inf}(\alpha _{\mathrm{ext}})>qt_{inf}(\alpha )=t_{inf}(\alpha ^q)\underset{s}{inf}(\alpha ^q).$$
By (i), every element of $`[\alpha ^q]^U`$ has a standard reduction system.
(vi) It can be proved in a way similar to (v). โ
Now, let us consider the following algorithm. Let $`\alpha `$ be a given non-periodic $`n`$-braid.
Applying cyclings and decyclings to $`\alpha `$, obtain an element $`\beta `$ of the set $`[\alpha ]^U[\alpha ]_๐^U`$ together with an element $`\gamma `$ such that $`\alpha =\gamma \beta \gamma ^1`$.
Decide whether $`\beta `$ has a standard reduction system or not.
If $`\beta `$ has no standard reduction system, then return โwe cannot decide whether $`\alpha `$ is reducibleโ, and halt.
Find a standard reduction system, say $`๐`$, of $`\beta `$.
Return โ$`\gamma ๐`$ is a reduction system of $`\alpha `$โ.
Note that, from definitions,
$$[\alpha ]^U[\alpha ]_๐^U=\{\beta [\alpha ]^S:๐^{\mathrm{}}(\beta )=\beta =๐^m(\beta )\text{for some }\mathrm{},m1\}.$$
This set is called the *reduced super summit set*, and known to be nonempty \[Lee00\].
Theorem 7.4 (i), (ii) and (iii) say that the above algorithm finds a reduction system of a non-periodic reducible braid $`\alpha `$ if either $`inf_s(\alpha _{\mathrm{ext}})>inf_s(\alpha )`$, $`sup_s(\alpha _{\mathrm{ext}})<sup_s(\alpha )`$, or $`\alpha `$ is a split braid. This implies that, roughly speaking, if the outermost component $`\alpha _{\mathrm{ext}}`$ is simpler than the whole braid $`\alpha `$ up to conjugacy, then we can find a reduction system of $`\alpha `$ from any element of $`[\alpha ]^U[\alpha ]_๐^U`$.
In Theorem 7.4, the conditions in (v) and (vi) are weaker than those in (i) and (ii). Because $`inf_s()`$ and $`sup_s()`$ are integer-valued, Theorem 2.9 (iii) implies the following.
* If $`inf_s(\alpha _{\mathrm{ext}})>inf_s(\alpha )`$, then $`inf_s(\alpha _{\mathrm{ext}})inf_s(\alpha )+1`$ and, hence,
$$t_{inf}(\alpha _{\mathrm{ext}})\underset{s}{inf}(\alpha _{\mathrm{ext}})\underset{s}{inf}(\alpha )+1>t_{inf}(\alpha ).$$
* If $`sup_s(\alpha _{\mathrm{ext}})<sup_s(\alpha )`$, then $`sup_s(\alpha _{\mathrm{ext}})sup_s(\alpha )1`$ and, hence,
$$t_{sup}(\alpha _{\mathrm{ext}})\underset{s}{sup}(\alpha _{\mathrm{ext}})\underset{s}{sup}(\alpha )1<t_{sup}(\alpha ).$$
Note that, for any $`m0`$, a braid $`\alpha `$ is reducible if and only if $`\alpha ^m`$ is reducible. Therefore, in order to decide the reducibility of $`\alpha `$, it suffices to decide the reducibility of $`\alpha ^m`$ for an arbitrary $`m0`$. If $`t_{inf}(\alpha _{\mathrm{ext}})>t_{inf}(\alpha )`$ or $`t_{sup}(\alpha _{\mathrm{ext}})<t_{sup}(\alpha )`$, then the above algorithm, applied to $`\alpha ^m`$ for $`1m<n(n1)/2`$, finds a reduction system of $`\alpha ^m`$ and, hence, decides the reducibility of $`\alpha `$. Consequently, the non-periodic reducible braids whose reducibility are not decidable by Theorem 7.4 are those with $`t_{inf}(\alpha _{\mathrm{ext}})=t_{inf}(\alpha )`$ and $`t_{sup}(\alpha _{\mathrm{ext}})=t_{sup}(\alpha )`$.
We close this section with some examples. From the examples, we can see that, in each statement of Theorem 7.4, the assertion does not hold if one of the conditions is weakened.
Example 7.5 shows that Theorem 7.4 (i), (ii) and (iii) do not hold for super summit sets. Namely, there is a split braid who satisfies the conditions (i) and (ii) but whose super summit set contains an element without standard reduction system.
###### Example 7.5.
Let $`\alpha =\sigma _1^1\sigma _2B_4`$ and $`\beta =(\sigma _3\sigma _2)^1\alpha (\sigma _3\sigma _2)=\sigma _2^1\sigma _1^1\sigma _2\sigma _3`$. (See Figure 14.) Then $`\alpha `$ is a split braid with
$$0=\underset{s}{inf}(\alpha _{\mathrm{ext}})>\underset{s}{inf}(\alpha )=1\text{and}0=\underset{s}{sup}(\alpha _{\mathrm{ext}})<\underset{s}{sup}(\alpha )=1$$
and $`\beta [\alpha ]^S`$, but $`\beta `$ has no standard reduction system.
Example 7.6 shows the following.
* Theorem 7.4 (i) and (ii) do not hold for $`inf_s(\alpha _{\mathrm{ext}})=inf_s(\alpha )`$ and $`sup_s(\alpha _{\mathrm{ext}})=sup_s(\alpha )`$, respectively. Namely, there is a non-periodic reducible braid $`\alpha `$ with $`inf_s(\alpha _{\mathrm{ext}})=inf_s(\alpha )`$ and $`sup_s(\alpha _{\mathrm{ext}})=sup_s(\alpha )`$ such that the set $`[\alpha ]^U[\alpha ]_๐^U`$ contains an element without standard reduction system.
* For a non-periodic reducible braid $`\alpha `$ with periodic $`\alpha _{\mathrm{ext}}`$, it is necessary to consider the ultra summit set $`[\alpha ^q]^U`$ of some power of $`\alpha `$ in Theorem 7.4 (iv). Namely, there is a non-periodic reducible $`\alpha `$ with periodic $`\alpha _{\mathrm{ext}}`$ such that $`[\alpha ]^U`$ contains an element without standard reduction system.
###### Example 7.6.
Consider the following 6-braids in Figure 15.
$`\alpha `$ $`=`$ $`\sigma _2\sigma _1\sigma _3\sigma _2\sigma _4\sigma _5\sigma _3\sigma _4\sigma _3`$
$`\beta `$ $`=`$ $`(\sigma _2\sigma _4^1)^1\alpha (\sigma _2\sigma _4^1)=\sigma _4\sigma _1\sigma _3\sigma _2\sigma _4\sigma _5\sigma _4\sigma _3\sigma _2`$
Observe that $`\alpha `$ is a non-periodic reducible braid such that $`\alpha _{\mathrm{ext}}=\sigma _1\sigma _2`$ is a periodic 3-braid. Since $`\alpha _{\mathrm{ext}}`$, $`\alpha `$ and $`\beta `$ are all permutation braids, we have
$$\underset{s}{inf}(\alpha )=0=\underset{s}{inf}(\alpha _{\mathrm{ext}});\underset{s}{sup}(\alpha )=1=\underset{s}{sup}(\alpha _{\mathrm{ext}});\beta [\alpha ]^U[\alpha ]_๐^U.$$
It is easy to see that $`\beta `$ has no standard reduction system.
Example 7.7 is due to Juan Gonzรกlez-Meneses and Bert Wiest. The authors are very grateful to them for providing it. It shows that Theorem 7.4 (v) and (vi) do not hold for $`t_{inf}(\alpha _{\mathrm{ext}})=t_{inf}(\alpha )`$ and $`t_{sup}(\alpha _{\mathrm{ext}})=t_{sup}(\alpha )`$, respectively. More precisely, there exist a non-periodic reducible braid $`\alpha `$ with $`t_{inf}(\alpha _{\mathrm{ext}})=t_{inf}(\alpha )`$ and $`t_{sup}(\alpha _{\mathrm{ext}})=t_{sup}(\alpha )`$, and an element $`\beta `$ such that, for each $`q1`$, the power $`\beta ^q`$ belongs to the set $`[\alpha ^q]^U[\alpha ^q]_๐^U`$ but has no standard reduction system.
###### Example 7.7.
Consider the following 7-braids in Figure 16.
$`\alpha `$ $`=`$ $`\sigma _1\sigma _2\sigma _3\sigma _4\sigma _3\sigma _2\sigma _1\sigma _5\sigma _4\sigma _6\sigma _5\sigma _4`$
$`\beta `$ $`=`$ $`(\sigma _3\sigma _4\sigma _5)^1\alpha (\sigma _3\sigma _4\sigma _5)=\sigma _1\sigma _2\sigma _3\sigma _2\sigma _1\sigma _4\sigma _3\sigma _5\sigma _6\sigma _5\sigma _4\sigma _3`$
Observe that
* both $`\alpha `$ and $`\beta `$ are permutation braids;
* $`\alpha `$ and $`\beta `$ are non-periodic reducible braids with reduction systems as in Figure 16;
* because $`\alpha _{\mathrm{ext}}`$ is pseudo-Anosov, the curves in Figure 16 (a) and (b) are the only reduction systems of $`\alpha ^q`$ and $`\beta ^q`$, respectively, for all $`q0`$.
Let $`B=\beta `$. (Throughout the paper, we have used capital letters $`A,B,\mathrm{}`$ to denote permutation braids.) The starting set and finishing set of $`B`$ are
$$S(B)=\{1,3,6\}\text{and}F(B)=\{1,3,4,6\}.$$
Since $`S(B)F(B)`$, the left normal form of $`\beta ^q`$ is $`\mathrm{\Delta }^0\underset{q}{\underset{}{BB\mathrm{}B}}`$ for all $`q1`$. In particular, for all $`q1`$,
$$๐(\beta ^q)=\beta ^q,๐(\beta ^q)=\beta ^q,inf(\beta ^q)=0\text{and}sup(\beta ^q)=q.$$
Therefore, for all $`q1`$, the power $`\beta ^q`$ belongs to the set $`[\alpha ^q]^U[\alpha ^q]_๐^U`$ and
$`t_{inf}(\alpha )`$ $`=`$ $`t_{inf}(\beta )=\underset{q\mathrm{}}{lim}inf(\beta ^q)/q=0;`$
$`t_{sup}(\alpha )`$ $`=`$ $`t_{sup}(\beta )=\underset{q\mathrm{}}{lim}sup(\beta ^q)/q=1.`$
The outermost component $`\alpha _{\mathrm{ext}}`$ is obtained from $`\alpha `$ by deleting the second strand. Similarly to the above, we can see that $`t_{inf}(\alpha _{\mathrm{ext}})=0=t_{inf}(\alpha )`$ and $`t_{sup}(\alpha _{\mathrm{ext}})=1=t_{sup}(\alpha )`$.
## 8. Proof of Proposition 7.3
In this section, we prove Proposition 7.3 that if $`\alpha `$ is a non-periodic reducible $`n`$-braid with $`inf_s(\alpha _{\mathrm{ext}})>inf_s(\alpha )`$, then for any element $`\beta `$ of $`[\alpha ]^U`$, the cycling commutator $`T_\beta `$ is a split braid.
Throughout this section, the notation $`\mathrm{St}^{\mathrm{ext}}(\gamma )`$ is used as an abbreviation for $`\mathrm{St}(_{\mathrm{ext}}(\gamma ))`$, the standardizer of the outermost component of the canonical reduction system of the braid $`\gamma `$. Therefore $`\mathrm{St}^{\mathrm{ext}}(\gamma )`$ consists of all positive braids $`P`$ such that $`P_{\mathrm{ext}}(\gamma )=_{\mathrm{ext}}(P\gamma P^1)`$ is standard. Recall that if $`\gamma [\gamma ]^U`$ and $`P`$ is the $`_R`$-minimal element of $`\mathrm{St}^{\mathrm{ext}}(\gamma )`$, then $`P\gamma P^1[\gamma ]^U`$ by Theorem 4.9.
Let $`\beta `$ be an element of the ultra summit set $`[\alpha ]^U`$. Then $`๐_0^m(\beta )=\beta `$ for some $`m1`$. For each $`i=0,\mathrm{},m`$, we define $`n`$-braids $`A_i`$, $`P_i`$ and $`\gamma ^{(i)}`$ as follows (see Figure 17):
* $`A_i`$ is the $`_R`$-minimal element of $`\{PB_n^+:inf(P๐_0^i(\beta ))>inf(๐_0^i(\beta ))\}`$;
* $`P_i`$ is the $`_R`$-minimal element of $`\mathrm{St}^{\mathrm{ext}}(๐_0^i(\beta ))`$;
* $`\gamma ^{(i)}=P_i๐_0^i(\beta )P_i^1`$.
Then, for each $`i=0,\mathrm{},m1`$,
* $`A_i`$ is a permutation braid with $`๐_0^{i+1}(\beta )=A_i๐_0^i(\beta )A_i^1`$ by Lemma 2.11 (i);
* $`_{\mathrm{ext}}(\gamma ^{(i)})`$ is standard because $`_{\mathrm{ext}}(\gamma ^{(i)})=_{\mathrm{ext}}(P_i๐_0^i(\beta )P_i^1)=P_i_{\mathrm{ext}}(๐_0^i(\beta ))`$ and $`P_i\mathrm{St}^{\mathrm{ext}}(๐_0^i(\beta ))`$;
* $`\gamma ^{(i)}`$ belongs to $`[\alpha ]^U`$ by Theorem 4.9.
###### Lemma 8.1.
For $`i=0,\mathrm{},m1`$, there exists a permutation braid $`B_i`$ such that $`B_iP_i=P_{i+1}A_i`$ and $`\gamma ^{(i+1)}=B_i\gamma ^{(i)}B_i^1`$.
###### Proof.
(See Figure 18.) Let $`B_i^{}`$ be the $`_R`$-minimal element of $`\{PB_n^+:inf(P\gamma ^{(i)})>inf(\gamma ^{(i)})\}`$. Then $`B_i^{}`$ is a permutation braid by Lemma 2.11, and
$$inf(B_i^{}\gamma ^{(i)})>inf(\gamma ^{(i)})\text{and}๐_0(\gamma ^{(i)})=B_i^{}\gamma ^{(i)}B_i^1.$$
Since both $`\gamma ^{(i)}`$ and $`๐_0^i(\beta )`$ belong to $`[\alpha ]^U`$, we have $`inf(\gamma ^{(i)})=inf(๐_0^i(\beta ))=inf_s(\alpha )`$. Since
$$inf(B_i^{}P_i๐_0^i(\beta ))=inf(B_i^{}\gamma ^{(i)}P_i)inf(B_i^{}\gamma ^{(i)})>inf(\gamma ^{(i)})=inf(๐_0^i(\beta )),$$
$`B_i^{}P_i`$ belongs to the set $`\{PB_n^+:inf(P๐_0^i(\beta ))>inf(๐_0^i(\beta ))\}`$. Since $`A_i`$ is the $`_R`$-minimal element of this set, we have $`A_i_RB_i^{}P_i`$, and hence
(7)
$$B_i^{}P_i=P_{i+1}^{}A_i$$
for some $`P_{i+1}^{}B_n^+`$. Note that
$$P_{i+1}^{}๐_0^{i+1}(\beta )P_{i+1}^1=P_{i+1}^{}A_i๐_0^i(\beta )A_i^1P_{i+1}^1=B_i^{}P_i๐_0^i(\beta )P_i^1B_i^1=B_i^{}\gamma ^{(i)}B_i^1=๐_0(\gamma ^{(i)}).$$
Since $`_{\mathrm{ext}}(๐_0(\gamma ^{(i)}))`$ is standard by Lemma 5.4, $`P_{i+1}^{}`$ belongs to $`\mathrm{St}^{\mathrm{ext}}(๐_0^{i+1}(\beta ))`$. Since $`P_{i+1}`$ is the $`_R`$-minimal element of $`\mathrm{St}^{\mathrm{ext}}(๐_0^{i+1}(\beta ))`$, we have $`P_{i+1}_RP_{i+1}^{}`$. Therefore,
(8)
$$P_{i+1}^{}=B_i^{\prime \prime }P_{i+1}$$
for some $`B_i^{\prime \prime }B_n^+`$. Observe that
$$P_{i+1}A_i๐_0^i(\beta )A_i^1P_{i+1}^1=P_{i+1}๐_0^{i+1}(\beta )P_{i+1}^1=\gamma ^{(i+1)}.$$
Since $`_{\mathrm{ext}}(\gamma ^{(i+1)})`$ is standard, $`P_{i+1}A_i`$ belongs to $`\mathrm{St}^{\mathrm{ext}}(๐_0^i(\beta ))`$. Since $`P_i`$ is the $`_R`$-minimal element of $`\mathrm{St}^{\mathrm{ext}}(๐_0^i(\beta ))`$, we have $`P_i_RP_{i+1}A_i`$. Therefore
(9)
$$P_{i+1}A_i=B_iP_i$$
for some $`B_iB_n^+`$. It is obvious that $`\gamma ^{(i+1)}=B_i\gamma ^{(i)}B_i^1`$. From (7), (8) and (9),
$$B_i^{}P_i=P_{i+1}^{}A_i=B_i^{\prime \prime }P_{i+1}A_i=B_i^{\prime \prime }B_iP_i.$$
Therefore $`B_i^{}=B_i^{\prime \prime }B_i`$. Since $`B_i^{}`$ is a permutation braid and $`B_i_RB_i^{}`$, the positive braid $`B_i`$ is a permutation braid as desired. โ
Let $`_{\mathrm{ext}}(\gamma ^{(0)})=๐_๐ง`$ for a composition $`๐ง=(n_1,\mathrm{},n_k)`$ of $`n`$. Let $`\mathrm{\Delta }_i`$ be the fundamental braid of $`B_{n_i}`$.
###### Lemma 8.2.
For $`i=0,\mathrm{},m1`$, $`_{\mathrm{ext}}(\gamma ^{(i)})=๐_๐ง`$ and $`B_i_R(\mathrm{\Delta }_1\mathrm{}\mathrm{\Delta }_k)`$.
###### Proof.
Using induction on $`i`$, it suffices to show the following:
> If $`_{\mathrm{ext}}(\gamma ^{(i)})=๐_๐ง`$, then $`B_i_R(\mathrm{\Delta }_1\mathrm{}\mathrm{\Delta }_k)`$ and $`_{\mathrm{ext}}(\gamma ^{(i+1)})=๐_๐ง`$.
Suppose $`_{\mathrm{ext}}(\gamma ^{(i)})=๐_๐ง`$. By Lemma 3.5 (ii) and (iv),
$$\gamma ^{(i)}=(\gamma _1\mathrm{}\gamma _k)\gamma _0_๐ง,$$
where $`\gamma _0=\gamma ^{(i)}{}_{\mathrm{ext}}{}^{}B_k`$ and $`\gamma _jB_{n_j}`$ for $`j=1,\mathrm{},k`$. Since $`inf_s(\alpha _{\mathrm{ext}})>inf_s(\alpha )`$ (from the hypothesis) and $`\gamma ^{(i)}[\alpha ]^U`$, we have $`inf(\gamma ^{(i)}{}_{\mathrm{ext}}{}^{})>inf(\gamma ^{(i)})`$ by Lemma 5.5. By Lemma 3.6,
$$inf(\gamma _0)=inf(\gamma ^{(i)}{}_{\mathrm{ext}}{}^{})>inf(\gamma ^{(i)})=\mathrm{min}\{inf(\gamma _i):i=0,\mathrm{},k,\mathrm{br}(\gamma _i)2\}.$$
Therefore $`inf(\gamma _0)>inf(\gamma _j)`$ for some $`j1`$ with $`\mathrm{br}(\gamma _j)2`$, and
$`inf((\mathrm{\Delta }_1\mathrm{}\mathrm{\Delta }_k)\gamma ^{(i)})=inf((\mathrm{\Delta }_1\gamma _1\mathrm{}\mathrm{\Delta }_k\gamma _k)\gamma _0_๐ง)`$
$`=`$ $`\mathrm{min}(\{inf(\mathrm{\Delta }_j\gamma _j):j=1,\mathrm{},k,\mathrm{br}(\gamma _j)2\}\{inf(\gamma _0)\})`$
$`=`$ $`\mathrm{min}(\{inf(\gamma _j)+1:j=1,\mathrm{},k,\mathrm{br}(\gamma _j)2\}\{inf(\gamma _0)\})`$
$`>`$ $`inf(\gamma ^{(i)}).`$
So $`(\mathrm{\Delta }_1\mathrm{}\mathrm{\Delta }_k)\{PB_n^+:inf(P\gamma ^{(i)})>inf(\gamma ^{(i)})\}`$. Recall, from the proof of Lemma 8.1, that $`B_i_RB_i^{}`$, where $`B_i^{}`$ is the $`_R`$-minimal element of $`\{PB_n^+:inf(P\gamma ^{(i)})>inf(\gamma ^{(i)})\}`$. Therefore,
$$B_i_RB_i^{}_R(\mathrm{\Delta }_1\mathrm{}\mathrm{\Delta }_k)$$
as desired. This implies that $`B_i`$ has the decomposition $`B_i=(B_{i,1}\mathrm{}B_{i,k})`$ for some permutation $`n_j`$-braid $`B_{i,j}`$โs. By Lemma 3.5 (ii), $`B_i๐_๐ง=๐_๐ง`$. Therefore, $`_{\mathrm{ext}}(\gamma ^{(i+1)})=_{\mathrm{ext}}(B_i\gamma ^{(i)}B_i^1)=B_i_{\mathrm{ext}}(\gamma ^{(i)})=B_i๐_๐ง=๐_๐ง`$. โ
Let $`S=B_{m1}\mathrm{}B_0`$. Then $`S`$ is a split braid by Lemma 8.2. Note that the cycling commutator of $`\beta `$ is $`T_\beta =A_{m1}\mathrm{}A_0`$. Since $`P_0^1SP_0=T_\beta `$ by Lemma 8.1, $`T_\beta `$ is a split braid and the proof is completed.
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# Novel order parameter to describe the critical behavior of Ising spin glass models
## 1 Introduction
Despite over three decades of intensive work, the nature of the low temperature phase of two-dimensional Edwards-Anderson (EA) model for spin glasses remains controversial. It is agreed that a phase transition occurs at zero temperature for a Gaussian distribution of bonds (GD) . Similarly, for a symmetric $`\pm J`$ distribution or bimodal distribution (BD) of bonds, very convincing numerical evidence has been found that there is no transition at finite temperature . In most of these references, the authors do not use an order parameter for characterizing the phase transition. On the other hand, data arising from other contributions, which are based on the behavior of a standard overlapping order parameter, support the existence of a finite critical temperature .
In this context, the main purposes of this paper are the following: a) To show that the disagreement pointed out in previous paragraph is related to the non-zero overlap of site-order parameters obtained for quite distinct energy valleys; b) To overcome this situation by proposing here a novel order parameter $`\mathrm{\Phi }`$, which is quite drastic to characterize phases but still is general enough to coincide with usual descriptions of ferromagnetic (F) and antiferromagnetic (AF) systems; c) To apply $`\mathrm{\Phi }`$ to do a scaling analysis for two-dimensional EA systems including Binder cumulant ; d) To confirm the assumption of the zero-temperature phase transition for two-dimensional BD, thus reinforcing this result obtained by previously quoted authors; and e) To give a physical meaning to this result by using the grounds on which $`\mathrm{\Phi }`$ is based on.
The present work is organized as it follows. In Section 2, we introduce the model and define a novel order parameter, $`\mathrm{\Phi }`$, very useful for spin glasses and other frustrated systems. Results of the simulation are presented in Section 3. Finally, our conclusions are drawn in Section 4.
## 2 Model and basic definitions
Let us begin by very briefly introducing the system under study. Ising spin $`s_i`$ occupies $`ith`$ site of a two dimensional (square for simplicity) lattice. The interaction with the spin at site $`j`$ is mediated by exchange interaction $`J_{ij}`$. In the absence of magnetic field (which is the case for the scope of the present paper) the Hamiltonian of such system can then be written as
$$H=\underset{i,j}{}J_{ij}s_is_j,$$
(1)
where interactions $`\{J_{ij}\}`$ are restricted to nearest neighbor couplings. In the ferromagnetic (F) Ising model, $`J_{ij}=J`$ $``$ $`i,j`$. For the EA model, we will consider half of the bonds F, while the other half will be described by antiferromagnetic (AF) bonds of the same magnitude, namely, $`J_{ij}=+J`$ ($`J>0`$). A sample is one of the possible random distributions of these mixed bonds. For simplicity spins take values $`s_j=\pm 1`$, which can be equally denoted by their signs.
Now, let us consider a configuration $`\alpha `$ defined by a collection of ordered spin orientations $`\{s_j^\alpha \}`$. The usual EA order parameter $`q`$ is built up by means of overlaps between two configurations $`\alpha `$ and $`\beta `$ and takes the form
$$q_{\alpha \beta }=\frac{1}{N}\underset{j=1}{\overset{N}{}}s_j^\alpha s_j^\beta ,$$
(2)
where $`N`$ ($`L\times L`$) is the total number of spins.
For models in which the ground state is non degenerate after breaking ergodicity, such as the pure F case, the distribution of $`q_{\alpha \beta }`$ values for the ground manifold (T=0.0) is trivial and it is given by delta functions at $`q_{\alpha \beta }=1.0`$ and $`q_{\alpha \beta }=1.0`$. This also happens in general for all systems with non-degenerate ground level. But this also applies to GD, where local fields have all different values at different sites, leading to a true minimum energy for just one pair of opposite ground states. However, for the BD the local field assumes a few discrete values only, which necessarily means highly degenerate ground manifolds leading to $`|q_{\alpha \beta }|<1.0`$, for a large number of possible pairs of ground states. This distribution will have two broad symmetric maxima but it will not vanish in the intermediate region .
On the other hand, a more detailed description based on a topological picture of the ground state of BD was presented . This framework allows us to define a state function with a clear physical meaning, which is a good candidate to be a new order parameter for a phase transition. In fact, it has been reported an important feature of the ground state, namely, at $`T=0`$ there exist clusters of solidary spins (CSS) preserving the magnetic memory of the system (solidary spins maintain their relative orientation for all states of the ground manifold). The main idea of this work is to characterize the nature of the low temperature phase through the CSS.
Let us consider a particular sample of any given size $`N`$. We denote by $`\mathrm{\Gamma }_\kappa `$ any of the $`n`$ CSS of the sample ($`\kappa `$ runs from $`1`$ to $`n`$). Calculations begin recognizing all of the CSS of each sample belonging to a set of $`2000`$ randomly generated samples of each size. This process is closely related to finding the so-called โdiluted latticeโ that prevails after removing all frustrated bonds , so the algorithms designed for that purpose can also be used here.
Let us first pick any arbitrary ground state configuration denoted by an asterisk ($``$) fixing one of the two possible relative orientations of the CSS, thus becoming a reference configuration. Then a local overlap corresponding to the configuration $`\alpha `$ in the $`\kappa `$-th cluster, of size $`N_\kappa `$, can be defined as
$$\varphi _\kappa ^\alpha =\frac{1}{N_\kappa }\underset{j\mathrm{\Gamma }_\kappa }{}s_j^{}s_j^\alpha ,$$
(3)
where the sum runs over all spins in the cluster $`\mathrm{\Gamma }_\kappa `$ only. Thus, $`|\varphi _\kappa ^\alpha |=1`$ indicates a fully ordered cluster; otherwise $`|\varphi _\kappa ^\alpha |<1`$. The magnetic order of the sample is characterized by the set of overlaps, namely, $`\{\varphi _\kappa ^\alpha \}`$. Under the occurrence of a phase transition, the new set $`\{\varphi _\kappa ^\alpha \}`$ will determine uniquely the ergodic component of the reached phase. This fact is a required characteristic for a well behaved order parameter .
We are now ready to define the new order parameter introduced in this paper. It is given by
$$\mathrm{\Phi }_\alpha =\underset{\kappa =1}{\overset{n}{}}f_\kappa |\varphi _\kappa ^\alpha |,$$
(4)
where $`f_\kappa =N_\kappa /N_I`$, being $`N_I=_{\kappa =1}^nN_\kappa `$, ($`N_IN`$). From the definition it flows that for $`TT_c`$ the average value of $`\mathrm{\Phi }_\alpha `$, namely, $`\mathrm{\Phi }`$, should be $`0`$. Similarly, for $`T<T_c`$ it should hold that $`0<\mathrm{\Phi }1`$, being $`\mathrm{\Phi }=1`$ for $`T=0`$ only. It is important to emphasize that $`\mathrm{\Phi }_\alpha `$ is a state function, which is an advantage over $`q_{\alpha \beta }`$ defined in eq. (2) as an overlap between two configurations of the system.
The calculation of the new order parameter $`\mathrm{\Phi }_\alpha `$ requires the previous determination of the set of CSS for each considered sample. This procedure, which was performed by using the numerical scheme introduced in Ref., is a computational limitation for going to larger system sizes. Once the ground manifold of each sample is completely characterized after this procedure, the numerical calculations converge very quickly by flipping the spins not present in the largest CSS only. The second run on each sample takes much less time than the first one that is needed to find all CSS.
In the F Ising model there is a unique cluster of $`N`$ solidary parallel spins at $`T=0`$. As it can be trivially demonstrated, eq. (4) leads to the magnetization per spin, which is the natural order parameter of such system. Similarly, for the AF case we get the well-known order parameter defined as the magnetization difference between the two possible interpenetrating sublattices. Finally, for GD, there also exists an unique CSS in the ground state and the phase transition occurring in the system is completely described by the new order parameter, eq.(4), as well. So the new parameter retains all the properties of the well-known non-degenerate systems.
Finally, the reduced fourth-order cumulant, introduced by Binder , can be calculated as
$$U_L=1\frac{\left[m^4\right]}{3\left[m^2\right]^2},$$
(5)
where $`m`$ is a given order parameter, and $`\mathrm{}`$ and $`[\mathrm{}]`$ mean the spin configuration (thermal) average and the bond configuration (sample) average, respectively. In general, the structure of the distribution of $`m`$ affects the behavior of the fourth-order Binder cumulant. Thus, for a trivial distribution, both $`|m|\pm 1`$ and $`U_L2/3`$, as $`T`$ goes to zero. On the other hand, if the distribution is nontrivial, $`|m|`$ tends to a value $`m^o`$ lower than 1, while $`U_L`$ tends to a value $`U_L^o`$ lower than $`2/3`$ upon decreasing temperature.
## 3 Results
Distribution functions for $`q_{\alpha \beta }`$ and $`\mathrm{\Phi }_\alpha `$, were obtained for BD by using a standard simulated-tempering procedure<sup>1</sup><sup>1</sup>1It must be emphasized that it is not necessary a simulated-tempering scheme for calculating the new order parameter, $`\mathrm{\Phi }`$. along with the well known Glauberโs dynamics . For illustration purposes, we perform calculations on $`1000`$ samples of size 64 $`(8\times 8)`$ at different temperatures ranging from $`T=0.2`$ to $`T=1.0`$. (throughout this paper, $`k_B/J=1`$ without any loss of generality). The results corresponding to $`T=0.31`$, $`T=0.53`$ and $`T=0.69`$ are presented in Fig. 1. As it is shown in part (a), the distribution of the new order parameter, $`R(\mathrm{\Phi }_\alpha )`$, exhibits a drastic behavior as $`T`$ decreases. In part (b) it is shown how the corresponding curves for $`r(|q_{\alpha \beta }|)`$ have a broad maximum over the plotted range and $`r(0)>0`$. These undesired characteristics for this order parameter remain even at low temperatures.
In Fig. 2, $`U_L(T)`$, built up from $`R(\mathrm{\Phi }_\alpha )`$ distribution, is presented for different lattice sizes ranging from $`N=16`$ to $`N=144`$ and each point was calculated by averaging over a set of 2000 samples. With the help of the inset, it is observed that the curves do not intersect each other as a direct indication of the absence of a phase transition for finite temperature, at least for the sizes considered here. Eventually we are not free from finite size considerations yet as it has been recently proposed that at least samples with $`L=50`$ should be reached when conventional parameters are used . However, using a more drastic parameter like the one proposed here, a faster convergence towards large $`L`$ values is expected. It is clear that all curves go to $`2/3`$ as $`T0`$, which reinforces the robustness of eq.(4). On the other hand, this property is not followed by cumulants obtained from other overlapping order parameters. This is the case of Fig. 7 in Ref., where it is possible to think that the reported crossing of the cumulants of $`q`$ arises from the dependence of $`U_L^o`$ on size. In this contribution, the authors reported a critical temperature different from zero, $`T_c0.23`$.
Finite-size scaling predicts that all curves in a figure such as Fig. 2, should collapse onto a single one when using $`(TT_c)L^{1/\nu }`$ as independent variable, being $`T_c`$ the critical temperature for the transition and $`\nu `$ an appropriate critical exponent. Upon choosing $`T_c=0`$ and the exponent $`\nu `$ is taken as $`\nu =2.63\pm 0.20`$, the standard universal behavior for $`U_L(T)`$ is obtained as shown in Fig. 3. This is an independent confirmation of previously reported results .
The following two parameters were also measured as each sample was solved exactly: (a) The mean fraction of spins, $`P\left[N_I\right]/N`$, belonging to any CSS; and (b) The fraction of spins in the largest CSS $`p\left[N_{\mathrm{}}\right]/N`$, where $`N_{\mathrm{}}`$ is the number of spins in the largest CSS. Fig. 4(a) shows that while $`P`$ remains rather constant, $`p`$ clearly decreases with size and the stabilizing role of the largest CSS is lost. The average number of CSS $`[n]`$ as function of size was also measured, finding that $`[n]`$ grows linearly with $`N`$, as it is shown in Fig. 4(b). For $`N>49`$, say (when small size effects do not play an important role), the following approximate law is obtained $`[n]0.03N+0.60`$.
The size dependence of a possible spin-glass phase can be described in the following terms. For small sizes ($`N<49`$ say) most of solidary spins are grouped in one large cluster stabilizing a spin-glass phase. As size grows, the number of CSS increases linearly with $`N`$ (or quadratically with $`L`$), while the relative size of the largest CSS diminishes. This can be visualized as if the original lattice would break into portions of relatively smaller sizes, none of them large enough to stabilize a spin-glass phase. This is the reason for the numeric result of Fig. 2, showing no intersection of curves for different sizes.
If the same procedure used here for the symmetric case is applied to different concentrations of F and AF bonds, a stable phase is found in the extremes of high and low concentrations of F bonds in correspondence with results already reported in the literature . As the relative concentration of F bonds varies the behavior of $`P`$ and $`[n]`$ is very similar to that shown in Fig. 4. However, $`p`$ tends to be constant for very asymmetric distributions of $`\pm J`$ bonds. The last statement indicates the presence of an infinite CSS in the thermodynamical limit, which is associated to a stable phase. These results are not shown graphically in the present paper.
## 4 Conclusions
A new order parameter $`\mathrm{\Phi }`$ has been introduced and applied to the study of magnetic systems. It proves to be particularly essential for characterizing degenerate systems such as Ising-like models with bimodal distribution. Parameter $`\mathrm{\Phi }`$ is well behaved, having all desired properties for a drastic order parameter. This behavior is based on the properties of CSS. When this order parameter is used for systems with BD, properties similar to order parameters for non-degenerate systems are found. Then, the characterization of magnetic phases after using the scaling techniques of cumulants becomes unambiguous. In this way it was shown that the two-dimensional Edwards-Anderson model exhibits a phase transition at $`T_c=0`$, with a critical exponent $`\nu =2.63\pm 0.20`$.
The identification of all CSS for each sample is the bottleneck in the present computational scheme. This procedure is very time consuming for large lattice sizes. The extra time needed for finding all CSS is well paid by the better precision achieved in the characterization of the phase, and the elimination of overlaps in the new order parameter, thus making the identification of the ergodic valley reliable.
The characteristics of this new order parameter make it also useful for other frustrated systems, where large overlaps occur due to the complex energy valley. The extension of the use of the parameter $`\mathrm{\Phi }`$ to other kind of problems is clearly foreseen. For instance, it can be the key element $`i)`$ to describe the phase diagram for the asymmetric distribution problem around the critical concentration of ferromagnetic (or antiferromagnetic) bonds of $`0.1`$ and $`ii)`$ to study the critical behavior of 3D Ising spin glasses. This task requires serious improvements in the numerical techniques used in order to get access to large lattice sizes. Work along this line is in progress.
Acknowledgments
We thank Fondecyt (Chile) under projects 1020993 and 7020993. One author (EEV) thanks Millennium Scientific Initiative (Chile) under contract P-02-054-F for partial support. Three authors (FR, FN and AJRP) thank CONICET (Argentina) and the Universidad Nacional de San Luis (Argentina) under project 322000.
## FIGURE CAPTIONS
Fig. 1. Distributions of the order parameters (a) $`\mathrm{\Phi }_\alpha `$ and (b) $`|q_{\alpha \beta }|`$, at 3 different temperatures as indicated.
Fig. 2. Cumulant $`U_L(T)`$ plotted versus $`T`$ for various lattice sizes as indicated. The inset zooms the area indicated by a dashed frame.
Fig. 3. Scaling plot of $`U_L`$ against $`(TT_c)L^{1/\nu }`$, with $`T_c=0`$ and $`\nu =2.63\pm 0.20`$.
Fig. 4. (a) $`P`$ and $`p`$ and (b) the growing of $`[n]`$ as a function of the lattice size, respectively.
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# On the accuracy of slow-roll inflation given current observational constraints
## I Introduction
Inflation is a theory which postulates that a rapid expansion of the universe occurred right after the Big Bang Guth (1981); Sato (1981); Albrecht and Steinhardt (1982); Linde (1982). Most inflationary models can be represented by an effective single field model with effective potential $`V`$. The inflaton with mass $`m`$ rolls down the potential until the kinetic energy of the inflaton is greater than half of its potential energy. At this point the inflationary expansion of the universe stops and the next phase of reheating occurs. During the inflationary expansion, the initial quantum fluctuations exponentially increase and become classical Starobinskii (1979); Mukhanov and Chibisov (1981); Guth and Pi (1982); Bardeen et al. (1983); Hawking (1982); Starobinsky (1982). These classical fluctuations also seed the subsequent growth of large scale structure. There is a well defined procedure which allows us to find the spectrum of the fluctuations given the inflationary potential. Because exact solutions are numerically intensive several appoximations have been developed. The most common approximation is the slow-roll approximation. Recently the so-called uniform approximation was suggested (Habib et al., 2004, 2005). Reference Casadio et al. (2005) developed improved WKB-type approximation.
If the kinetic energy of the inflaton is much smaller than its potential energy, we say that the inflaton is slowly rolling down its potential. In this slow-roll approximation we can obtain analytical formulae for the produced power spectrum in the form of a Taylor series expansion in a set of slow-roll parameters. The coefficients in the Taylor expansion of the logarithm of the power spectrum in $`\mathrm{ln}k`$ effectively define the slope $`n_s1`$, running $`\alpha _s`$ and higher derivatives. We usually derive the slow-roll formulae through the time delay formalism or Bessel function approximation. Therefore there are some implied conditions on the accuracy of the slow-roll approximation depending upon the slow-roll parameters. References Wang et al. (1997); Leach et al. (2002) found that there are areas in the slow-roll parameter space where the accuracy of slow-roll approximation is questionable. This usually requires a large deviation of $`n_s`$ from 1. However it contradicts the latest observations (Spergel et al., 2003; Verde et al., 2003; Seljak et al., 2005).
Recently there has been a lot of renewed interest in models with large running of the scalar index Peiris et al. (2003); Kawasaki et al. (2003); Chung et al. (2003). It is not clear whether slow-roll approximation is accurate in this area of parameter space, as in some expansions one of the slow roll parameters becomes large and the expansion is no longer well controlled Leach et al. (2002). Another issue is the question of where to stop the expansion. Although it is often assumed that the running is $`O\left((n_s1)^2\right)`$, reference Dodelson and Stewart (2002) found that there are cases where it can be as large as $`n_s1`$. In this case one should also consider the effect of including the running of the running of $`n_s`$, i.e. the second derivative of $`n_s`$ over $`\mathrm{ln}k`$. These are the issues addressed in this paper. We begin with a short review of the basic physics of inflation and the algorithm of numerical solutions to the inflationary equations, with more details given in appendix. We continue by comparing the numerical solutions to those given by slow-roll approximations and finally we present our conclusions.
In this paper we use a standard convention for reduced Planck mass $`m_{\mathrm{pl}}=G_N^{1/2}`$.
## II Inflationary basics
In the โHamilton-Jacobiโ formulation, the evolution of the Hubble parameter $`H(\varphi )`$ during inflation with potential $`V(\varphi )`$ is given by (e.g. see Liddle and Lyth (2000))
$$\left[H^{}(\varphi )\right]^2\frac{12\pi }{m_{\mathrm{pl}}^2}H^2(\varphi )=\frac{32\pi ^2}{m_{\mathrm{pl}}^4}V(\varphi ).$$
(1)
The number of e-folds $`N`$ since some initial time is related to the value of the scalar field $`\varphi `$ by
$$\frac{dN}{d\varphi }=\frac{4\pi }{m_{\mathrm{pl}}^2}\frac{H(\varphi )}{H^{}(\varphi )}.$$
(2)
We will consider the situation when the value of the scalar field is growing in time, $`d\varphi /dt>0`$. Then by our convention $`dN/dt`$ is also positive, $`dN/dt>0`$.
In the literature, different sets of slow-roll parameters are used. Reference Liddle and Lyth (1992) introduce potential slow-roll parameters which are constructed on the basis of the derivatives of the inflationary potential $`V(\varphi )`$. Authors of Liddle et al. (1994) define Hubble slow-roll parameters through the derivatives of the Hubble parameter $`H(\varphi )`$ with respect to the field $`\varphi `$ during inflation
$`ฯต_H(\varphi )`$ $`=`$ $`{\displaystyle \frac{m_{\mathrm{pl}}^2}{4\pi }}\left({\displaystyle \frac{H^{}(\varphi )}{H(\varphi )}}\right)^2,`$ (3)
$`\eta _H(\varphi )`$ $`=`$ $`{\displaystyle \frac{m_{\mathrm{pl}}^2}{4\pi }}{\displaystyle \frac{H^{\prime \prime }(\varphi )}{H(\varphi )}},`$ (4)
$`{}_{}{}^{n}\xi _{H}^{}(\varphi )`$ $`=`$ $`\left({\displaystyle \frac{m_{\mathrm{pl}}^2}{4\pi }}\right)^n{\displaystyle \frac{(H^{})^{n1}H^{(n+1)}}{H^n}}.`$ (5)
In this parameterization when the inequality $`ฯต_H(\varphi )<1`$ fails, the inflation immediately stops. Sometimes $`{}_{}{}^{2}\xi _{H}^{}`$ is also denoted $`\xi _H`$ or $`\xi _H^2`$ though it can take negative values. In this paper we will use $`\xi _H{}_{}{}^{2}\xi _{H}^{}`$.
Reference Schwarz et al. (2001) introduces another basis of โhorizon-flowโ slow-roll parameters through the logarithmic derivative of the Hubble distance $`ฯต_0=d_H=1/H(N)`$ with respect to the number of e-folds $`N`$ to the end of inflation
$$ฯต_{n+1}=\frac{d\mathrm{ln}|ฯต_n|}{dN}.$$
(6)
The connection between any two of these sets can be found in e.g. Schwarz et al. (2001). Thus the first three horizon-flow slow-roll parameters are connected to the first three Hubble slow-roll parameters as (Liddle et al., 1994; Lidsey et al., 1997; Leach et al., 2002)
$`ฯต_1`$ $`=`$ $`ฯต_H,`$ (7)
$`ฯต_2`$ $`=`$ $`2ฯต_H2\eta _H,`$ (8)
$`ฯต_2ฯต_3`$ $`=`$ $`4ฯต_H^26ฯต_H\eta _H+2\xi _H.`$ (9)
There is an analytical connection between Hubble slow-roll parameters and potential slow-roll parameters (Liddle et al., 1994).
References Martin and Schwarz (2000); Stewart and Lyth (1993); Stewart and Gong (2001); Wei et al. (2004) use differently defined sets of slow-roll parameters, but they still can be converted to the ones we have described here (e.g. see (Schwarz et al., 2001)).
Thus any inflationary model can be completely described by the evolution of one of the sets of the parameters.
The condition for the inflation to occur is $`ฯต_1=ฯต_H<1`$ or $`ฯต_V1`$ since $`ฯต_H=ฯต_V`$ to first order.
To find the power spectrum of the perturbations produced by a single field inflation, one can follow the prescription of Grivell and Liddle (Grivell and Liddle, 1996). One solves the equation (Mukhanov, 1985, 1988; Stewart and Lyth, 1993)
$$\frac{d^2u_k}{d\tau ^2}+\left(k^2\frac{1}{z}\frac{d^2z}{d\tau ^2}\right)u_k=0$$
(10)
for each mode with wavenumber $`k`$ and initial condition $`u_k(\tau )\frac{1}{\sqrt{2k}}e^{ik\tau }`$ as $`\tau \mathrm{}`$. Then the spectrum of curvature perturbations is given by
$$๐ซ_{}(k)=\frac{k^3}{2\pi ^2}\left|\frac{u_k}{z}\right|^2.$$
(11)
The quantity $`z`$ in equation (10) is defined as $`z=a\dot{\varphi }/H`$ for scalar modes and $`z=a`$ for tensor modes. Then for scalar modes (Grivell and Liddle, 1996)
$$\begin{array}{c}\frac{1}{z}\frac{d^2z}{d\tau ^2}=2a^2H^2[1+ฯต_H\frac{3}{2}\eta _H\hfill \\ \hfill +ฯต_H^22ฯต_H\eta _H+\frac{1}{2}\eta _H^2+\frac{1}{2}\xi _H].\end{array}$$
(12)
One can parametrize the power spectrum of the scalar and tensor modes of the fluctuations amplified by the inflation as
$$\mathrm{ln}\frac{๐ซ(k)}{๐ซ_0}=(n1)\mathrm{ln}\frac{k}{k_{}}+\frac{\alpha }{2}\mathrm{ln}^2\frac{k}{k_{}}+\frac{\beta }{6}\mathrm{ln}^3\frac{k}{k_{}}+\mathrm{}$$
(13)
around some conventional pivot point $`k_{}`$. Leach et al. (Leach et al., 2002) give expressions for the scalar spectral index $`n_s`$, the running of the scalar spectral index $`\alpha _s`$, the tensor spectral index $`n_t`$ and the running of the tensor spectral index $`\alpha _t`$ in terms of the horizon-flow parameters. Here we reproduce their second order formulae for $`n_s1`$ and $`\alpha _s`$
$`n_s1`$ $`=`$ $`2ฯต_1ฯต_22ฯต_1^2`$ (14)
$`(2C+3)ฯต_1ฯต_2Cฯต_2ฯต_3,`$
$`\alpha _s`$ $`=`$ $`2ฯต_1ฯต_2ฯต_2ฯต_3,`$ (15)
where $`C=\gamma _\mathrm{E}+\mathrm{ln}220.7296`$.
Reference Leach et al. (2002) also analyzes the accuracy of the approximation (13) for parameterizing the inflationary power spectrum of fluctuations with $`\beta =0`$ for different values of the parameters $`r`$, $`n_s`$ and $`\alpha _s`$.
In this paper we will also compare the second order formulae (14, 15) to the first order formulae given by
$`n_s1`$ $`=`$ $`2ฯต_1ฯต_2,`$ (16)
$`\alpha _s`$ $`=`$ $`2ฯต_1ฯต_2ฯต_2ฯต_3.`$ (17)
The expression for $`\alpha _s`$ is the same as in the second order formula because the expression for $`\alpha _s`$ is derived using only first order expression for $`n_s`$. Thus the main difference between first and second order formulae comes from the extra terms in the expression for $`n_s1`$.
Current observational constraints on $`r`$, $`n_s`$ and $`\alpha _s`$ are given by Tegmark et al. (2004); Seljak et al. (2005); Slosar et al. (2004). At 95% confidence level, the tensor to scalar ratio is $`R<0.50`$, which implies that the first horizon-flow parameter $`ฯต_1`$ is much smaller than one. Current constraints on the scalar spectral index give us $`n_s=0.98\pm 0.02`$, which in turn means that the second horizon-flow slow-roll parameter is much smaller than one.
Present data does not require the presence of running in the primordial power spectrum (Liddle, 2004), but running as large as $`\pm 0.03`$ is still allowed at 3-$`\sigma `$ (Seljak et al., 2005). Regular inflationary models usually predict $`|\alpha _s|(n_s1)^2`$ and so the running is of the order of $`10^3`$, as is the case for the minimally-coupled $`V(\varphi )=\lambda \varphi ^4`$ model with 60 e-folds remaining.
But it is possible that $`|\alpha _s|(n_s1)^2`$ and $`\alpha _s<0`$, which means that the main part in the running of the spectral index (15) is determined not by the first term $`2ฯต_1ฯต_2`$, but by the second term $`ฯต_2ฯต_3`$. It happens when $`|ฯต_3||ฯต_1|`$, and therefore there might be a situation when $`|ฯต_3|>1`$.
To summarize, if $`n_s1`$ and $`\alpha _s`$ is a small negative number, at some scale we might have $`ฯต_11`$, $`ฯต_21`$ and $`|ฯต_3|>1`$. Leach et al. (Leach et al., 2002) define inflation satisfying slow-roll under the condition $`|ฯต_n|1`$, for all $`n>0`$. In our case $`ฯต_3>1`$, so the question arises as to whether slow-roll in this case is accurate or whether the approximation breaks down and one must also include terms with higher powers in $`ฯต_3`$. Does it mean that the inflation is not slow-roll and one must use full numerical solutions instead? And does it mean that one must also include the running of the running? These are the main questions we address in this paper. To address them we have developed the numerical code described in appendix A.
How natural is it for inflation with a given number from 50 to 70 e-folds remaining to produce a power spectrum with a changing tilt? In the absence of theoretical guidance on the inflationary space we cannot address this question simply. Authors of Peiris et al. (2003) have produced about 200,000 simulations of the inflationary flow equations for more or less โrandomโ potentials, and calculated the observable parameters ($`n_s`$, $`\alpha _s`$, $`r`$, $`n_t`$, $`\alpha _t`$) of the resulting power spectra about 40 to 70 e-folds before the end of the inflation for each potential. About 80,000 of them fall into the area plotted on Fig. 1. Only the fifteen marked with larger yellow circles give a significant change in the tilt from red to blue, i.e. $`n_s1`$, $`\alpha _s<0.02`$.
Choosing the Hubble parameter to be represented by a Taylor expansion in $`\varphi `$ with uniformly distributed coefficients, as done in Peiris et al. (2003), does not necessarily correspond to the real inflationary priors (Liddle, 2003). We do not address this issue here; instead we want to simply stress that possibility of constructing a potential with a large running in the scalar power spectrum 40-70 e-folds before the end of the inflation exists.
## III Quadratic potential
As a test of our code, in this section we investigate how well the slow-roll formulae work in the slow-roll regime for one of the usual potentials that do not predict large running. As an example we will consider a simple quadratic potential, $`V=m^2\varphi ^2/2`$, which is a classic example of chaotic inflation. The second panel from the bottom in Fig. 2 effectively shows the dependence of $`z^{\prime \prime }/z`$ on the number of the e-folds for inflation with such a potential. The behavior is monotonic and very smooth, which is due to the smoothness of the derivatives of the potential. Since $`z^{\prime \prime }/z`$ scales as $`2a^2H^2`$, we plot the quantity
$$\frac{1}{2a^2H^2}\frac{z^{\prime \prime }}{z}1$$
(18)
instead (compare to equation (12)).
The top two panels show the dependence of $`ฯต_H`$, $`\eta _H`$ and $`\xi _H`$ on the number of e-folds. The only significantly non-zero term is $`ฯต_H`$, which gradually grows to $`1`$ at the end of inflation. The values of $`\eta _H`$ and $`\xi _H`$ are typically smaller by roughly $`10^3`$ and $`10^4`$ respectively.
Figure 3 shows the primordial power spectrum produced by the quadratic potential. The second panel from the bottom describes the error produced by the slow-roll approximations. The first order approximation gives less than 0.2% error in the observed range of wavenumbers $`k`$. The second order approximation works slightly better; the error is just above 0.1%. Both of these numbers are likely to be good enough for the upcoming experiments. This is because the accuracy at large scales is limited by the finite number of modes, while at small scales it is limited by the nonlinear evolution. So, while the overall amplitude could in principle be determined to an accuracy of 0.1% when CMB and lensing information is combined, it is unlikely that such a precision will be achieved separately at two widely separated length scales.
Taking a more careful look at the error plot, one sees that the error curve in the observed area is basically a straight line, meaning that the main source of error is not the imprecise value of $`\alpha _s`$ but the error in $`n_s`$. Let us estimate now how precisely we need to know $`n_s`$ to get an error of, say 0.2%, in the observed range. The imprecision $`\delta n_s`$ in $`n_s`$ will give us the uncertainty
$$\delta n_s\frac{1}{2}\mathrm{ln}\frac{k_{\mathrm{max}}}{k_{\mathrm{min}}}=0.002.$$
(19)
Taking the observed range of $`k`$โs to be from $`10^3`$ Mpc<sup>-1</sup> to $`1`$ Mpc<sup>-1</sup>, we find that one needs to find $`n_s`$ with the precision of $`\delta n_s=610^4`$.
The same allowed uncertainty $`\delta \alpha _s`$ in $`\alpha _s`$ is estimated from
$$\frac{1}{2}\delta \alpha _s\left(\frac{1}{2}\mathrm{ln}\frac{k_{\mathrm{max}}}{k_{\mathrm{min}}}\right)^2=0.002.$$
(20)
Therefore $`\delta \alpha _s=310^4`$ is the error which we can make in determining $`\alpha _s`$ in order to get an error in the power spectrum of 0.2% at the edges of the observed range of $`k`$โs.
The two top panels of Fig. 3 compare the numerically found dependence of $`n_s`$ and $`\alpha _s`$ on $`k`$ to the one found from the slow-roll approximation with the first order $`\alpha _s`$ and either the first or second order for $`n_s`$. One should compare the discrepancies between these to the values of $`\delta n_s`$ and $`\delta \alpha _s`$. The characteristic value of $`n_s`$ is .964 and the discrepancy between the exact value and the one found from the slow-roll approximation is comparable to $`\delta n_s`$. Running $`\alpha _s`$ takes values around $`6.510^4`$. The discrepancy between the exact and the slow-roll values is very small in comparison to $`\delta \alpha _s`$. One can also notice that in this case $`|\alpha _s|2\delta \alpha _s`$. Thus, even if we assigned $`\alpha _s=0`$, we would not get a significant error in the approximation of the primordial power spectrum of the scalar perturbations.
To summarize this section, for standard inflationary potentials, the slow-roll approximation suffices even at first order when compared to the expected accuracy of existing and future experiments. The second order approximation, while improving the accuracy, is not really necessary. The main error of slow-roll when considered in contrast to the numerical solutions is the inaccuracy in the slope $`n_s`$; inaccuracies in higher order expansion terms, such as the running, are less important and can even be ignored.
## IV Potential with a bump in the second derivative
We want to construct a potential which will give us a strong running and crossing of the point $`n_s=1`$ in the observable power spectrum. We want to have $`n_s>1`$ at earlier times in inflation, while at later times we want to have $`n_s<1`$. To get the desired result, one can take two different potentials producing such features and smoothly connect them.
One can rewrite slow-roll formulae (14,15) through the potential slow-roll parameters as
$`n_s1`$ $`=`$ $`6ฯต_V+2\eta _V,`$ (21)
$`\alpha _s`$ $`=`$ $`16ฯต_V\eta _V24ฯต_V^22\xi _V.`$ (22)
Now let us just choose our potential to be
$$V(\varphi )=10.01\varphi 1.20\varphi ^2$$
(23)
for all $`\varphi >0`$. This choice provides about 50 e-folds of inflation after $`\varphi =0`$. Since the local properties of the power spectrum are mostly determined by the local โhistoryโ of the slow-roll parameters at the moment of the horizon crossing, we can get a red tilt of the scalar power spectrum $`n_s0.80`$ in the area where the โhistoryโ before point $`\varphi =0`$ is not very important. To get an approximately symmetric shape of the power spectrum we choose $`V(\varphi )`$ to be
$$V(\varphi )=10.01\varphi +1.20\varphi ^2$$
(24)
for all $`\varphi <0`$. In this case for wave modes which cross the horizon far before the moment when the scalar field takes the value of $`\varphi =0`$, the spectral index of the primordial power spectrum has a blue tilt $`n_s1.20`$. Thus between these two regions the spectral index changes from 1.20 to 0.80. We can unite formulae (23) and (24) into
$$V(\varphi )=10.01\varphi 1.20\varphi ^2sign\varphi .$$
(25)
This potential has continuous first and second derivatives, but has a bump in its third derivative. This makes $`{\displaystyle \frac{1}{z}}{\displaystyle \frac{d^2z}{d\tau ^2}}`$ in (12) discontinuous around $`\varphi =0`$. According to Starobinskii (1992) this produces oscillations in the power spectrum, which we can indeed see for the potential (25). To avoid the oscillations we smooth out the $`sign\varphi `$ function, changing it to $`{\displaystyle \frac{2}{\pi }}\mathrm{arctan}(200\varphi )`$. In this case (25) changes to
$$V(\varphi )=10.01\varphi 1.20\varphi ^2\frac{2}{\pi }\mathrm{arctan}(200\varphi ),$$
(26)
which is shown on Fig. 4. We have chosen 200 as the coefficient in front of $`\varphi `$ in the $`\mathrm{arctan}`$ function so that the produced power spectrum has a nice shape as in Fig. 5.
The two top panels of Fig. 4 show the behavior of the slow-roll parameters $`ฯต_H`$, $`\eta _H`$, $`\xi _H`$ and $`ฯต_1`$, $`ฯต_2`$, $`ฯต_3`$, $`ฯต_2ฯต_3`$ correspondingly. While nothing unexpected happens to the behavior of the conventional Hubble slow-roll parameters $`ฯต_H`$, $`\eta _H`$ and $`\xi _H`$, there appears to be a singularity for the horizon-flow parameter $`ฯต_3`$. However, notice that the product $`ฯต_2ฯต_3`$ behaves smoothly and remains small due to the fact that the parameter $`ฯต_2`$ is changing its sign and therefore crossing through zero. Thus the parameterization of equation (6) introduces a singularity which is not physically present in the model.
Figure 5 shows the power spectrum produced by the model of inflation with the potential (25). The second panel from the bottom shows the errors made by different approximations. We again observe a similar picture for the slow-roll formulae. The main source of error for either the first or second order approximations comes not from the value of $`\alpha _s`$ but from the error in the value of $`n_s`$. From the second panel from the top, we can estimate that the discrepancy is of the order of 0.01 for
$$\delta n_s=n_s^{\mathrm{exact}}n_s^{\mathrm{approx}}$$
(27)
which gives an error of
$$\delta n_s\frac{1}{2}\mathrm{ln}\frac{k_{\mathrm{max}}}{k_{\mathrm{min}}}4\%$$
(28)
in the produced power spectrum at the edges of the observed range. Both the first and second order slow-roll approximations for $`n_s`$ work somewhat unsatisfactory. The first order slow roll underestimates $`n_s`$ and the second order overestimates it by about the same amount.
On the other hand, if our goal is to focus on running alone regardless of the slope and use just that property to deduce something about the potential, then the slow-roll does very well, since the differences between the slow-roll and numerical value of running are very small even at the lowest order in slow-roll. Extra terms in the expansion (13) further improve the accuracy. Adding running of the running improves the accuracy over the observed range from 1% to 0.2%.
In summary, for potentials that lead to large running, slow-roll does not estimate the slope $`n_s`$ very accurately at either first or second order, while the accuracy of the running $`\alpha _s`$ suffices for the existing and future experiments. If we observe over a wide range of scales then it is useful to add the cubic term. Second order slow-roll does not seem to improve the accuracy.
## V Flow Equations Simulations
Kinney Kinney (2002) introduced a formalism based on the so-called flow equations, further discussed in Liddle (2003). The basic idea is that if one fixes the Hubble slow-roll parameters (5) at some point in time for $`ฯต_H`$, $`\eta _H`$ and $`{}_{}{}^{\mathrm{}}\xi _{H}^{}`$ up to $`\mathrm{}=M`$ and assumes that all the other Hubble slow-roll parameters are small enough that one can neglect them in oneโs calculations (i.e. $`{}_{}{}^{\mathrm{}}\xi _{H}^{}=0`$ for all $`\mathrm{}M+1`$) then, without any other assumptions about inflation being slow-roll, one can find the Hubble slow-roll parameters at any other moment of time using the following hierarchy of linear ordinary differential equations:
$`{\displaystyle \frac{dฯต}{dN}}`$ $`=`$ $`2ฯต(\eta ฯต),`$ (29)
$`{\displaystyle \frac{d\eta }{dN}}`$ $`=`$ $`ฯต\eta {}_{}{}^{2}\xi ,`$ (30)
$`{\displaystyle \frac{d{}_{}{}^{\mathrm{}}\xi }{dN}}`$ $`=`$ $`[\mathrm{}ฯต(\mathrm{}1)\eta ]{}_{}{}^{\mathrm{}}\xi {}_{}{}^{\mathrm{}+1}\xi `$ (31)
for all $`\mathrm{}=2\mathrm{}M`$ assuming $`{}_{}{}^{M+1}\xi =0`$.
Usually when we set up an inflationary problem, we choose a potential $`V(\varphi )`$ and then reconstruct the form of the Hubble parameter during inflation using the main nonperturbed Hamilton-Jacobi inflationary equation (1), which gives us an attractor solution $`H(\varphi )`$ which in the inflationary class of problems almost does not depend on the initial condition.
By following the method prescribed by Kinney (2002) one avoids solving the main attractor inflationary equation (1), as pointed out by Liddle (2003). Indeed, the assumption $`{}_{}{}^{\mathrm{}}\xi _{H}^{}=0`$ for all $`\mathrm{}M+1`$ requires that $`H^{(\mathrm{})}(\varphi )=0`$ for all $`\mathrm{}M+2`$. Consequently, $`H(\varphi )`$ is a polynomial of order $`M+1`$:
$$H(\varphi )=H_0(1+A_1\varphi +A_2\varphi ^2+A_3\varphi ^3+\mathrm{}+A_{M+1}\varphi ^{M+1}).$$
(32)
In this case the function $`H(\varphi )`$ is an attractor solution of equation (1) with a potential in the form
$$\begin{array}{c}V(\varphi )=\frac{m_{\mathrm{pl}}^4}{32\pi ^2}\left([H^{}(\varphi )]^2\frac{12\pi }{m_{\mathrm{pl}}^2}H^2(\varphi )\right)\hfill \\ \hfill =\frac{m_{\mathrm{pl}}^4}{32\pi ^2}H_0^2[(A_1+\mathrm{}+(M+1)A_{M+1}\varphi ^M)^2\\ \hfill \frac{12\pi }{m_{\mathrm{pl}}^2}(1+A_1\varphi +\mathrm{}+A_{M+1}\varphi ^{M+1})^2].\end{array}$$
(33)
Thus the only differential equation one needs to solve in order to match up the number of e-folds and the value of the scalar field $`\varphi `$ is
$$\frac{dN}{d\varphi }=\frac{2\sqrt{\pi }}{m_{\mathrm{pl}}}\frac{1}{\sqrt{ฯต(\varphi )}}.$$
(34)
Here again $`ฯต(\varphi )`$ is defined as in the equation (5):
$$\begin{array}{c}ฯต(\varphi )=\frac{m_{\mathrm{pl}}^2}{4\pi }\left(\frac{H^{}(\varphi )}{H(\varphi )}\right)^2\hfill \\ \hfill =\frac{m_{\mathrm{pl}}^2}{4\pi }\left(\frac{A_1+2A_2\varphi +\mathrm{}+(M+1)A_{M+1}\varphi ^M}{1+A_1\varphi +A_2\varphi ^2+\mathrm{}+A_{M+1}\varphi ^{M+1}}\right)^2.\end{array}$$
(35)
We do not have to numerically solve the hierarchy of $`M`$ differential flow equations. Instead we have analytical expressions for $`V(\varphi )`$ and $`H(\varphi )`$.
The late attractor $`ฯต={}_{}{}^{\mathrm{}}\xi =0`$ and $`\eta =const`$, found by Kinney (2002), corresponds to the situation where the inflation proceeds to the value of the scalar field $`\varphi `$, which is a solution of the equation $`ฯต=0`$
$$A_1+2A_2\varphi +\mathrm{}+(M+1)A_{M+1}\varphi ^M=0.$$
(36)
At this point if $`A_20`$, then $`\eta =const0`$ due to the definition of $`\eta `$, which does not involve $`ฯต`$ at all. All the other $`{}_{}{}^{\mathrm{}}\xi =0`$ since any of them is a product of the first derivative of the Hubble parameter (which is zero) with some higher order derivatives.
Peiris et al. Peiris et al. (2003) made $`M=9`$-th order flow equation simulations; about $`40,000`$ are shown as black dots on Fig. 1. Fifty points fall into the range $`|n_s1|<0.05`$ and $`\alpha _s<0.02`$; these are shown in yellow. Among these point we have chosen 13 which fall into the narrow interval $`|n_s1|<0.02`$, and we have reconstructed the corresponding inflationary potentials for the inflationary models which give such significant running, together with $`n_s`$ extremely close to $`1`$.
The bottom panel in Fig. 6 shows a potential from such a model with an unusually high value of the running $`\alpha _s`$. We notice that there is a small dip in the potential. Some of the potentials with high $`\alpha _s`$ from the simulations had unrealistically high values of the tensor to the scalar ratio, but all of them had quite similar shapes. The second from the bottom panel of Fig. 6 shows the characteristic behavior of the function $`z^{\prime \prime }/z`$ which influences the scalar power spectrum as we have seen earlier. The top two panels show the dependence of the slow-roll parameters on the number of e-folds. As in the other case with large running, we find a singularity for the horizon-flow slow-roll parameter $`ฯต_3`$, while the product $`ฯต_2ฯต_3`$ behaves smoothly and $`ฯต_2`$ crosses zero.
The bottom panel in Fig. 7 shows the power spectrum of scalar and tensor perturbations produced by inflation with the potential under consideration. The second from the bottom panel shows the error produced by every one of the approximations for the power spectrum. We again see that both first and second order slow-roll formulae do not give a satisfactory result for $`n_s`$. One of them again overestimates $`n_s`$; the other underestimates it. The error for either of the approximations is about 2โ4%.
The error introduced by the approximate formula for the running $`\alpha _s`$ is a bit smaller than the one for $`n_s`$, but it is somewhat larger compared to the quadratic potential we considered in the previous section.
Chen et al. (Chen et al., 2004) perform a similar analysis of slow-roll approximation. Using the flow-equations technique, they found discrepancy of larger than $`0.01`$ for $`n_s`$ between second and third order slow-roll approximations for some of the models. Based on this fact they conclude that third order slow-roll is better. For the model we considered in this section we have found that the third order slow-roll does not improve the results of the second order approximation. In our calculations both formulas give identical results leading to approximately the same order of error as the first order approximation.
## VI Is Truncated Taylor Expansion good?
Recently Abazajian, Kadota and Stewart (Abazajian et al., 2005) have argued that if
$$|\alpha _s\mathrm{ln}(k/k_{})||n_s1|,$$
(37)
then the traditional truncated Taylor series parameterization is inconsistent, and hence it can lead to incorrect parameter estimations. One can notice that Taylor expansions $`P(x)=a_ix^i`$ of functions $`x^2`$ or $`\mathrm{cos}x`$ around $`x=0`$ also violates the condition $`a_1a_2x`$, but no one argues that these expansions are not valid. Abazajian et al. propose to use the parameterization
$$\mathrm{ln}๐ซ(k)=\mathrm{ln}๐ซ_0+\frac{(n_s1)^2}{\alpha _s}\left[\left(\frac{k}{k_{}}\right)^{\frac{\alpha _s}{n_s1}}1\right]$$
(38)
instead.
There is one significant disadvantage of this approach. In particular, using this parameterization to describe $`๐ซ`$ as a function of $`k`$, one is able to describe only a growing or decreasing function, which can be easily seen from the form of the function. The models we study in this paper produce scalar power spectra which are not purely growing or decreasing (e.g. see Fig. 7).
In the previous section we have considered the potential which produces power spectrum satisfying equation (37). On Fig. 8 we compare the traditional truncated to second and third order Taylor expansion and the parameterization (38). We find that the parameterization (38) gives a significantly larger error than, e.g. the second order Taylor expansion.
Thus we found that in this particular case though equation (37) holds, truncated Taylor expansion is a good approximation and the AKS approach does not improve it. There might be models for which equation (38) works better than Taylor expansion, but it is definitely not an improvement for a general case and should be used with caution, if at all.
## VII Conclusions
In this paper we have explored the accuracy of the slow-roll approximation given the observational constraints on the primordial scalar and tensor power spectra. The current constraints can be roughly described by the tensor to scalar ratio $`r<1`$, small deviation from the scale invariance of the scalar power spectrum, $`|n_s1|<0.05`$ and small but possibly nontrivial running, $`|\alpha _s|<0.03`$. These constraints allow for the particular case where $`n_s1`$ and $`|\alpha _s|>0.01`$, which has previously been argued to not satisfy the slow-roll condition. We have computed exact numerical solutions for the considered potentials and compared them to those obtained from the first and second order slow-roll approximations.
We have found that for the potentials explored here, there is no substantial difference when using first or second order slow-roll formulae for the power spectrum index $`n_s`$. Both of them either work well in the case of small running or have a comparable error in the case of non-negligible running. Adding extra (cubic in $`\mathrm{ln}k`$) terms in the approximation for the scalar power spectrum extends the accuracy to a larger range of scales, but this accuracy is most likely not necessary for existing and near future experiments. If the values of $`n_s`$ and $`\alpha _s`$ are known with the precision $`\delta n_s=610^4`$ and $`\delta \alpha _s=310^4`$, then the scalar power spectrum will have an error of about 0.2% at the edge of the observable range of wavenumbers $`k`$โs.
The horizon-flow basis $`ฯต_{n+1}=d\mathrm{ln}|ฯต_n|/dN`$ introduces an artificial singularity for inflationary models with negative running and the value of the spectral index crossing 1. Such a divergence in one of the horizon-flow parameters does not indicate that the slow-roll approximation has been badly broken. We find that the slow-roll is still accurate at the 1-2% level and most of the error comes from inaccuracies in the evaluation of the slope itself, and not the running. Thus the first order slow-roll approximation is sufficiently accurate for the current observations. Only if the running turns out to be large, while the slope remains close to scale-invariant, are exact numerical calculations required to achieve sub-percent accuracy. In the appendix we present a short guideline on performing such calculations. One can request the code directly from the author.
###### Acknowledgements.
The author is thankful to Hiranya Peiris for making the results of her Monte Carlo simulations available. AM also thanks Sergei Bashinsky, Chris Beasley, Latham Boyle, Steven Gratton, Patricia Li, Uroลก Seljak and Alexei Starobinskii for useful discussions and comments.
*
## Appendix A Inflationary Equations
In this appendix we describe the technical details of the code we ran to get the results presented in the main part of the paper. The code is given a potential $`V(\varphi )`$ and some point $`\varphi _0`$ which lies in the observable range of wave-modes and, say, corresponds to the moment when wavelengths with $`k=0.05`$ Mpc<sup>-1</sup> exit the horizon. We want to find the power spectrum produced by inflation with the potential $`V(\varphi )`$. For this purpose we first have to go backwards in time about 50 e-folds and then start the inflation there. This guarantees that the inflationary dynamics are not affected by the choice of the initial condition and we indeed have the attractor solution.
Now we evolve the universe from our โbeginning of inflationโ to the end of inflation, the moment which is determined by the violation of the inequality $`\ddot{a}>0`$. This part is described below in the โnon-perturbed inflationary equationsโ section. Usually we require 50 to 70 e-folds between $`\varphi _0`$ and the end of inflation.
After we already have the complete background history of the evolution of the universe during the inflationary stage of the expansion, we can start working out the evolution of the perturbations during inflation, as discussed in the second part of the appendix.
### A.1 Non-perturbed inflationary equations
The unperturbed dynamics of inflation are described by the equation of motion of the scalar field $`\varphi `$ with potential $`V(\varphi )`$ in the expanding universe with the Hubble parameter $`H\dot{a}/a`$
$$\ddot{\varphi }+3H\dot{\varphi }+V^{}(\varphi )=0$$
(39)
and the Friedman equation with only the scalar field component present in the universe
$$H^2=\frac{8\pi }{3m_{\mathrm{pl}}^2}\left[V(\varphi )+\frac{1}{2}\dot{\varphi }^2\right].$$
(40)
The equations (3940) are equivalent to the pair of Hamilton-Jacobi equation (1) and
$$\dot{\varphi }=\frac{m_{\mathrm{pl}}^2}{4\pi }H^{}(\varphi ).$$
(41)
The Hamilton-Jacobi equation connects the Hubble parameter and the value of the potential of the scalar field during the inflation. In the case when we know the behavior of the Hubble parameter it is easy to find the potential. The method of flow equations is entirely based on this fact. In contrast, if we know the shape of the potential and want to reconstruct the behavior of the Hubble parameter, the problem is not as simple. First of all, as for any first order differential equation, we would like to have an initial condition $`H_0=H(\varphi _0)`$. Due to the attractor nature of the equation (1) its solution does not really depend on the initial condition $`H_0`$ (we have found from numerical simulations that one needs about 6 e-folds to forget the history). Thus it does not really matter which initial condition we choose.
Hamilton-Jacobi equation requires that
$$H^2(\varphi )\frac{8\pi }{3m_{\mathrm{pl}}^2}V(\varphi ).$$
(42)
If we are going to use a method such as Runge-Kutta for the integration of the differential equation(1), we might try values of $`H`$ which would violate the inequality (42).
To avoid this complication, we reparametrize our equation using a new function $`\delta (\varphi )`$ so that
$$H^2(\varphi )=\frac{8\pi }{3m_{\mathrm{pl}}^2}V(\varphi )\left(1+e^{\delta (\varphi )}\right).$$
(43)
Then substituting our new definition into equation (1) we get
$$H^{}(\varphi )=\frac{4\pi \sqrt{2}}{m_{\mathrm{pl}}^2}\sqrt{V(\varphi )}e^{\delta (\varphi )/2}.$$
(44)
Combining this with the expression for $`H^{}`$ obtained from the direct differentiation of $`H`$ in (43), we get a differential equation for $`\delta ^{}(\varphi )`$
$$\delta ^{}=\sqrt{1+e^\delta }\left[\frac{V^{}}{V}\sqrt{1+e^\delta }+\frac{4\sqrt{3\pi }}{m_{\mathrm{pl}}}\right].$$
(45)
This equation is much more pleasant to deal with numerically than equation (1), since it does not have a weird boundary for $`\delta `$, as $`H`$ did before. One can also check the attractor nature of the equation (45), that it does not remember the prior history. We see now that in the case when the potential is changing slowly $`\delta ^{}0`$ and we have
$$e^\delta =\left[\frac{48\pi }{m_{\mathrm{pl}}^2\left(V^{}/V\right)^2}1\right]^1\frac{m_{\mathrm{pl}}^2}{48\pi }\left(\frac{V^{}}{V}\right)^2.$$
(46)
We can use this approximate solution of the equation as the initial condition for our differential equation since it is quite close to the true solution and it will make our numerical solution evolve into the attractor solution faster.
One can check that the expressions for $`ฯต_H`$, $`\eta _H`$ and $`{}_{}{}^{2}\xi _{H}^{}`$ are given by the following formulae
$`ฯต`$ $`=`$ $`{\displaystyle \frac{m_{\mathrm{pl}}^2}{4\pi }}\left({\displaystyle \frac{H^{}}{H}}\right)^2={\displaystyle \frac{3}{1+e^\delta }},`$ (47)
$`\eta `$ $`=`$ $`{\displaystyle \frac{m_{\mathrm{pl}}^2}{4\pi }}{\displaystyle \frac{H^{\prime \prime }}{H}}`$ (48)
$`=`$ $`3+{\displaystyle \frac{m_{\mathrm{pl}}}{4}}\sqrt{{\displaystyle \frac{3}{\pi }}}{\displaystyle \frac{V^{}}{V}}{\displaystyle \frac{1}{\sqrt{e^\delta (1+e^\delta )}}},`$
$`{}_{}{}^{2}\xi `$ $`=`$ $`{\displaystyle \frac{m_{\mathrm{pl}}^4}{16\pi ^2}}{\displaystyle \frac{H^{}H^{\prime \prime \prime }}{H^2}}`$ (49)
$`=`$ $`3(ฯต+\eta )\eta ^2{\displaystyle \frac{3m_{\mathrm{pl}}^2}{8\pi }}{\displaystyle \frac{V^{\prime \prime }}{V}}{\displaystyle \frac{1}{1+e^\delta }}.`$
From these expressions we can expect that in general $`ฯต`$ and $`\eta `$ are continuous functions, whereas $`{}_{}{}^{2}\xi `$ does not have to be continuous at points where $`V^{\prime \prime }`$ is not continuous.
The condition for inflation to take place ($`\ddot{a}>0`$) follows from the derivative of the Friedman equation
$$\frac{\ddot{a}}{a}=\frac{8\pi }{3m_{\mathrm{pl}}^2}\left[V(\varphi )\dot{\varphi }^2\right]=H^2(\varphi )(1ฯต),$$
(50)
or
$$\begin{array}{c}\frac{\ddot{a}}{a}=\frac{8\pi }{3m_{\mathrm{pl}}^2}\left[V(\varphi )\frac{m_{\mathrm{pl}}^4}{16\pi ^2}\left(H^{}(\varphi )\right)^2\right]\hfill \\ \hfill =\frac{8\pi }{3m_{\mathrm{pl}}^2}V(\varphi )(12e^\delta ).\end{array}$$
(51)
The first of these two equations implies that the end of inflation happens when the inequality $`ฯต<1`$ is violated. The same thing occurs when the inequality $`\delta <\mathrm{ln}2`$ is violated in the second equation. The latter also means that the inflation continues while the kinetic energy of the inflaton is less than half of its potential energy
$$\frac{K}{\mathrm{\Pi }}=\frac{\dot{\varphi }^2/2}{V(\varphi )}=e^{\delta (\varphi )}<\frac{1}{2}.$$
(52)
Thus we come to a physical definition of our parameter $`e^{\delta (\varphi )}`$ as the ratio of the kinetic energy $`\dot{\varphi }^2/2`$ to the potential energy $`V(\varphi )`$.
In the next subsection we will be working with inflationary perturbations and it will not be very convenient for us to work with the value of the scalar field $`\varphi `$ as an independent variable. For this purpose we will use the number of e-folds defined as
$$\stackrel{~}{N}=\mathrm{ln}\frac{(aH)}{(aH)_0}.$$
(53)
Note that this is the actual number of e-folds and is not the same as $`N=\mathrm{ln}(a/a_0)`$. The connection between $`\stackrel{~}{N}`$ and $`\varphi `$ is determined through the derivative
$$\frac{d\stackrel{~}{N}}{d\varphi }=\frac{2\sqrt{\pi }}{m_{\mathrm{pl}}}\frac{1ฯต(\varphi )}{\sqrt{ฯต(\varphi )}}.$$
(54)
We are almost done describing the background evolution of the universe, except we have not yet chosen the initial value of the scalar field $`\varphi _i`$. We only have the value $`\varphi _0`$ which corresponds to the moment when the mode $`k=0.05`$ Mpc<sup>-1</sup> exits the horizon. We want to move backwards in time for about 50 e-folds. Equation (45) has an attractor behavior only when we are moving in the positive direction along the $`\varphi `$-axis. It diverges from the attractor solution in the negative direction. As a useful trick, let us modify equation (1) to the following form
$$[H^{}(\varphi )]^2=\frac{12\pi }{m_{\mathrm{pl}}^2}\left[\frac{8\pi }{3m_{\mathrm{pl}}^2}V(\varphi )H^2(\varphi )\right].$$
(55)
In this form, when we move backwards in time the value of $`H(\varphi )`$ is bound by the value of $`\sqrt{8\pi V/3m_{\mathrm{pl}}^2}`$ from the top and the solution cannot diverge. In addition we temporarily redefine $`\delta (\varphi )`$ to satisfy
$$H^2(\varphi )=\frac{8\pi }{3m_{\mathrm{pl}}^2}V(\varphi )(1e^{\delta (\varphi )}).$$
(56)
Thus we get an equation analogous to the equation (45)
$$\delta ^{}=\sqrt{1e^\delta }\left[\frac{V^{}}{V}\sqrt{1e^\delta }+\frac{4\sqrt{3\pi }}{m_{\mathrm{pl}}}\right].$$
(57)
Equation (57) does not carry any physical meaning; we just use this equation to go โupwardsโ to the higher values of the potential, still tracking the general behavior of $`V(\varphi )`$. If we go backwards in time 50 e-folds using (57) and then forward in time 50 e-folds using (45), we will not return to the same point $`\varphi _0`$, since the behavior of $`\delta (\varphi )`$ in the equation (57) is determined by the area which is to the right of the current value of $`\varphi `$ and in the equation (45) is determined by the area which is on the left side. Nevertheless, this method gives us a good estimate of what initial value of $`\varphi _i`$ we should take.
It is also worth mentioning that this approach is not more difficult to deal with than the inflationary flow equations (2931).
### A.2 Perturbation equations
#### A.2.1 Scalar mode
The algorithm for finding the scalar mode primordial power spectrum is described in the main text (see equation (10) and below). Here we will just mention some technical details.
Equation (10) is not very convenient to solve in its current form. First of all we would like to set the independent variable, the conformal time $`\tau `$, in such a way that $`\tau 0`$ as inflation goes on. In this case we would be able to numerically integrate equation (10) up to as small values of $`\tau `$ as we want. But in numerical realizations we cannot really choose such an initial value of $`\tau _i`$ that gives us $`\tau 0`$ at the end of the inflation. Suppose that at the end of the inflation we have $`\tau 10`$. In this case the numerical error on $`\tau `$ will be of the order of $`10^{15}`$ which is a reasonable machine precision. Hence the limit on corresponding $`d\tau `$ is of the same order and we can explore the range of changing the scale factor $`a`$ from $`1`$ to $`10^{15}`$, i.e. about 35 e-folds. This might be enough, but to be safe we will use a different independent variable, the true number of e-folds $`\stackrel{~}{N}`$ defined by equation (53) which is the same as
$$d\stackrel{~}{N}=\frac{d(aH)}{aH}.$$
(58)
Then the mode equation (10) can be rewritten as
$$\begin{array}{c}(1a)\frac{d^2u_k}{d\stackrel{~}{N}^2}+(1+b)\frac{du_k}{d\stackrel{~}{N}}\hfill \\ \hfill +\left[\left(\frac{k}{k_0}\right)^2e^{2(\stackrel{~}{N}\stackrel{~}{N}_0)}2(1+c)\right]u_k=0,\end{array}$$
(59)
where coefficient $`a`$, $`b`$ and $`c`$ can be exactly expressed through $`ฯต_H`$, $`\eta _H`$ and $`{}_{}{}^{2}\xi _{H}^{}`$ as
$`a`$ $`=`$ $`2ฯตฯต^2,`$ (60)
$`b`$ $`=`$ $`2ฯตฯต^2+2ฯต\eta ,`$ (61)
$`c`$ $`=`$ $`ฯต{\displaystyle \frac{3}{2}}\eta +ฯต^22ฯต\eta +{\displaystyle \frac{1}{2}}\eta ^2+{\displaystyle \frac{1}{2}}{}_{}{}^{2}\xi .`$ (62)
In equation (59), $`k`$ is the wavelength of the interest, while $`k_0`$ and $`\stackrel{~}{N}_0`$ are constants conveniently chosen for normalization purposes.
Further, equation (10) has a solution
$$u_k\frac{1}{\sqrt{2k}}e^{ik\tau }$$
(63)
at the beginning of the inflation when $`\tau \mathrm{}`$ and $`k^2{\displaystyle \frac{1}{z}}{\displaystyle \frac{d^2z}{d\tau ^2}}`$. We also know the approximate behavior of $`u_k`$ at later times when $`\tau 0`$ and $`k^2{\displaystyle \frac{1}{z}}{\displaystyle \frac{d^2z}{d\tau ^2}}`$:
$$u_kz.$$
(64)
Thus it is natural to decompose $`u_k`$ into growing and oscillating parts
$$u_k=e^{A+i\varphi },$$
(65)
where both functions $`A`$ and $`\varphi `$ are real functions of conformal time $`\tau `$ or of the true number of e-folds $`\stackrel{~}{N}`$. Then the equation (59) can be split into 4 ordinary differential equations with 2 new functions $`A_p`$ and $`\varphi _p`$ defined as below
$`{\displaystyle \frac{dA}{d\stackrel{~}{N}}}`$ $`=`$ $`A_p,`$ (66)
$`{\displaystyle \frac{dA_p}{d\stackrel{~}{N}}}`$ $`=`$ $`{\displaystyle \frac{(k^2/k_0^2)e^{2\stackrel{~}{N}}+2(1+c)(1+b)A_p}{1a}}`$ (67)
$`(A_p^2\varphi _p^2),`$
$`{\displaystyle \frac{d\varphi }{d\stackrel{~}{N}}}`$ $`=`$ $`\varphi _p,`$ (68)
$`{\displaystyle \frac{d\varphi _p}{d\stackrel{~}{N}}}`$ $`=`$ $`\varphi _p{\displaystyle \frac{(1+b)+2(1a)A_p}{1a}}.`$ (69)
This system of differential equations looks a bit more complicated than the single equation (10), but it is actually much easier to solve numerically. Indeed, at earlier times we have $`dA/d\stackrel{~}{N}A_p=0`$. This instaneously gives us the initial condition on $`d\varphi /d\stackrel{~}{N}\varphi _p`$ from (67)
$$\varphi _p^2=\frac{(k^2/k_0^2)e^{2\stackrel{~}{N}}2(1+c)}{1a}$$
(70)
as $`\tau \mathrm{}`$, i.e. $`\stackrel{~}{N}\mathrm{}`$. To be consistent with the initial condition on
$$u_k\frac{1}{\sqrt{2k}}e^{ik\tau }$$
as $`\tau \mathrm{}`$ we also require that
$$A=\frac{1}{2}\mathrm{ln}k$$
as $`\stackrel{~}{N}\mathrm{}`$. As the inflation continues, the terms
$$\frac{k^2}{k_0^2}\frac{e^{2\stackrel{~}{N}}}{1a}$$
and $`\varphi _p^2`$ will balance each other on the right hand side of the equation (67) until $`A_p`$ is not negligible in comparison to $`1`$ in equation (69). Thus, around $`\stackrel{~}{N}=\mathrm{ln}(k/k_0)`$ the oscillating part $`\varphi _p`$ will decrease more rapidly than before, finally exponentially dropping to zero. At the same time $`A_p`$, and therefore $`A`$, start exponentially growing. The final power spectrum is given by
$$๐ซ_k=\frac{k^3}{2\pi ^2}\left|\frac{u_k}{z}\right|^2\frac{k^3}{2\pi ^2}e^{2A_k}.$$
(71)
Thus we even do not need information about the phase $`\varphi `$ and we can freely drop equation (68) from our system. Also while being in the stage of inflation where $`u_k`$ has an oscillatory behavior, if one were to use the usual method without our substitution, one would have to find the values of $`u_k`$ for at least 6 points per oscillation period. However with our substitution, we easily pass this area, which does not have any interest for us since we analytically know the behavior of $`u_k`$ here, and therefore move directly to the place where we cannot solve it analytically. By our estimates this technique gives a gain of a factor of 10 in computational time, which is of particular interest if one wants to calculate the power spectrum for e.g. 100 wavemodes.
#### A.2.2 Tensor mode
The calculation of the tensor mode power spectrum of perturbations is absolutely analogous to the one for scalars, except instead of equation (10) one has to solve
$$\frac{d^2u_k}{d\tau ^2}+\left(k^2\frac{1}{a}\frac{d^2a}{d\tau ^2}\right)u_k=0$$
(72)
with the same initial condition
$$u_k(\tau )\frac{1}{\sqrt{2k}}e^{ik\tau }$$
as $`\tau \mathrm{}`$, where $`a`$ is the usual scale factor of the Friedman universe. One can show that
$$\frac{1}{a}\frac{d^2a}{d\tau ^2}=2a^2H^2\left(1\frac{1}{2}ฯต\right).$$
(73)
Mode equations for the amplitude and the phase of the wave (66-69) of the tensor mode look similar except in the equations (60-62) where we have to change $`c`$ to $`d`$ defined as
$$d=\frac{1}{2}ฯต.$$
(74)
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# Quasi-morphismes et invariant de Calabi
## 1. Introduction
Un quasi-morphisme sur un groupe $`\mathrm{\Gamma }`$ est une fonction $`\varphi :\mathrm{\Gamma }`$ telle que les quantitรฉs
$$\varphi (xy)\varphi (x)\varphi (y)$$
soient bornรฉes lorsque $`x,y`$ dรฉcrivent $`\mathrm{\Gamma }`$. On appellera parfois dรฉfaut de $`\varphi `$ la quantitรฉ
$$\delta =\mathrm{sup}_{x,y\mathrm{\Gamma }}|\varphi (xy)\varphi (x)\varphi (y)|.$$
Un quasi-morphisme est homogรจne sโil satisfait en outre $`\varphi (x^n)=n\varphi (x)`$ pour $`x\mathrm{\Gamma }`$ et $`n`$. Nous dirons que deux quasi-morphismes sont รฉquivalents si leur diffรฉrence est bornรฉe. Il nโest pas difficile de vรฉrifier que si $`\varphi `$ est un quasi-morphisme quelconque, la formule
$$\varphi _h(x)=\mathrm{lim}_p\mathrm{}\frac{1}{p}\varphi (x^p)$$
dรฉfinit lโunique quasi-morphisme homogรจne ร distance bornรฉe de $`\varphi `$. Il satisfait
$$|\varphi \varphi _h|\delta .$$
On pourra consulter (par exemple) pour une introduction ร ce sujet. Nous noterons $`QM_h(\mathrm{\Gamma },)`$ lโespace des quasi-morphismes homogรจnes sur le groupe $`\mathrm{\Gamma }`$.
Si $`(V,\omega )`$ est une variรฉtรฉ symplectique connexe fermรฉe, le groupe $`\mathrm{Ham}(V,\omega )`$ de ses diffรฉomorphismes hamiltoniens est simple dโaprรจs un thรฉorรจme de A. Banyaga . Il nโadmet donc pas de morphisme non-trivial vers $``$. Si $`(V,\omega )`$ est une variรฉtรฉ symplectique connexe ouverte, sur laquelle $`\omega `$ est exacte, E. Calabi a introduit dans un morphisme
$$๐๐ฉ_V:\mathrm{Ham}(V,\omega ).$$
Le noyau de ce morphisme est simple dโaprรจs un autre thรฉorรจme de Banyaga . Si $`\lambda `$ est une primitive de $`\omega `$ sur $`V`$ et $`(f_t)`$ une isotopie hamiltonienne dans $`V`$, engendrรฉe par le champ de vecteurs $`Z_t`$, on a :
$$๐๐ฉ_V(f_1)=_V_0^1\lambda (Z_t)๐t\omega .$$
Supposons maintenant que $`(V,\omega )`$ est fermรฉe. A chaque ouvert connexe $`UV`$, on associe le sous-groupe $`\mathrm{\Gamma }_U`$ de $`\mathrm{Ham}(V,\omega )`$ formรฉ des diffรฉomorphismes qui sont le temps $`1`$ dโune isotopie hamiltonienne dans $`U`$. Si $`\omega `$ est exacte sur $`U`$, on dispose alors du morphisme de Calabi : $`๐๐ฉ_U:\mathrm{\Gamma }_U`$. On notera $`๐`$ la famille des ouverts connexes $`U`$ de $`V`$ tels que $`\omega `$ est exacte sur $`U`$ et tels quโil existe $`f\mathrm{Ham}(V,\omega )`$ avec $`f(U)\overline{U}=\mathrm{}`$. Dans , M. Entov et L. Polterovich posent la question suivante :
Peut-on construire un quasi-morphisme homogรจne $`\varphi :\mathrm{Ham}(V,\omega )`$ dont les restrictions aux sous-groupes $`(\mathrm{\Gamma }_U)_{U๐}`$ soient รฉgales aux morphismes de Calabi $`(๐๐ฉ_U)_{U๐}`$ ?
Plus gรฉnรฉralement, peut-on construire un tel quasi-morphisme sur le revรชtement universel $`\stackrel{~}{\mathrm{Ham}}(V,\omega )`$ du groupe des diffรฉomorphismes hamiltoniens ? Dans (voir aussi ), ils rรฉpondent positivement ร cette question pour une certaine classe de variรฉtรฉs symplectiques, qui inclut notamment les espaces projectifs complexes, en particulier la sphรจre $`S^2`$. Leur mรฉthode utilise des outils sophistiquรฉs de topologie symplectique.
Dans lโesprit de constructions prรฉcรฉdentes de J.-M. Gambaudo et ร. Ghys, nous construisons un quasi-morphisme ayant une propriรฉtรฉ comparable ร celle รฉnoncรฉe dans la question ci-dessus, sur le groupe des diffรฉomorphismes hamiltoniens dโune surface de genre supรฉrieur ou รฉgal ร $`2`$. Dans , on pourra trouver de nombreuses constructions de quasi-morphismes sur les groupes $`\mathrm{Ham}(S,\omega )`$, pour toute surface fermรฉe orientรฉe (les quasi-morphismes de Gambaudo et Ghys sont en fait dรฉfinis sur le groupe $`\mathrm{Symp}_0(S,\omega )`$). Cependant les restrictions de ces quasi-morphismes aux sous-groupes $`(\mathrm{\Gamma }_U)`$ ne sont pas des homomorphismes. Les constructions de Gambaudo et Ghys peuvent รชtre vues comme de possibles gรฉnรฉralisations du nombre de translation $`\tau `$ sur le groupe $`\stackrel{~}{\mathrm{Hom}\stackrel{ยด}{\mathrm{e}}\mathrm{o}}_+(S^1)`$ des homรฉomorphismes de la droite rรฉelle qui commutent aux translations entiรจres. Elles sont dans lโesprit de constructions prรฉcรฉdentes par V.I. Arnold , D. Ruelle , S. Schwartzman .
###### Thรฉorรจme 1.
Soit $`S`$ une surface fermรฉe orientรฉe de genre supรฉrieur ou รฉgal ร $`2`$, munie dโune forme symplectique $`\omega `$. Il existe un quasi-morphisme homogรจne
$$๐๐ฉ_S:\mathrm{Ham}(S,\omega ),$$
dont la restriction aux sous-groupes $`\mathrm{\Gamma }_U`$ est รฉgale au morphisme de Calabi, dรจs que $`U`$ est diffรฉomorphe ร un disque ou ร un anneau. Le quasi-morphisme $`๐๐ฉ_S`$ est invariant par conjugaison par tout diffรฉomorphisme symplectique.
Nous pouvons faire deux remarques concernant ce thรฉorรจme. Dโune part, en utilisant des constructions de Gambaudo et Ghys, on peut le renforcer en lโรฉnoncรฉ :
lโespace (affine) des quasi-morphismes homogรจnes sur $`\mathrm{Ham}(S,\omega )`$ ayant la propriรฉtรฉ du thรฉorรจme $`1`$ est de dimension infinie.
Cela rรฉsultera simplement du fait que lโespace des quasi-morphismes homogรจnes sur $`\mathrm{Ham}(S,\omega )`$ dont les restrictions aux groupes $`\mathrm{\Gamma }_U`$ (oรน $`U`$ est diffรฉomorphe ร un disque ou ร un anneau) sont nulles, est un espace de dimension infinie. Dโautre part, la nature de ce quasi-morphisme est certainement trรจs diffรฉrente de celle du quasi-morphisme construit par Entov et Polterovich dans . En effet, parmi les deux conditions qui dรฉfinissent la famille dโouverts $`๐`$, la premiรจre est vide en dimension $`2`$, en revanche la seconde nโapparaรฎt pas du tout dans notre travail. Il existe dโailleurs des disques (ou des anneaux) ne la satisfaisant pas.
Lโรฉnoncรฉ du thรฉorรจme suivant est inspirรฉ du thรฉorรจme $`5.2`$ de , oรน un calcul similaire est fait sur le groupe $`\mathrm{Ham}(S^2,\omega )`$. Notre mรฉthode est cependant diffรฉrente. Considรฉrons donc une fonction de Morse $`F:S`$, dont les valeurs critiques sont toutes distinctes. Notons $`x_1,\mathrm{},x_l,`$ ses points critiques, $`\lambda _j=F(x_j)`$ ses valeurs critiques, avec $`\lambda _1<\mathrm{}<\lambda _l`$. On a classiquement :
$$\underset{j=1}{\overset{l}{}}(1)^{\mathrm{ind}x_j}=22g,$$
$`g`$ est le genre de $`S`$. Considรฉrons lโespace
$$=\{H:S,\{H,F\}=0\},$$
des fonctions sur la surface $`S`$ qui commutent avec $`F`$ au sens de Poisson, cโest-ร -dire telles que $`\omega (X_H,X_F)=0`$, oรน $`X_G`$ dรฉsigne le gradient symplectique dโune fonction $`G`$. Lโensemble
$$\mathrm{\Gamma }=\{\phi _H^1,H\}$$
est un sous-groupe abรฉlien de $`\mathrm{Ham}(S,\omega )`$ (oรน $`\phi _H^t`$ dรฉsigne le flot de $`X_H`$). La restriction de $`๐๐ฉ_S`$ ร $`\mathrm{\Gamma }`$ est donc un homomorphisme, que nous calculons dans le thรฉorรจme suivant. La donnรฉe de la fonction $`F`$ permet de construire un graphe fini $`๐ข`$ appelรฉ graphe de Reeb , de la maniรจre suivante. Parmi les composantes connexes des niveaux $`F^1(cste)`$ on trouve :
1. les points critiques de $`F`$ dโindice $`0`$ ou $`2`$,
2. des courbes simples plongรฉes,
3. des courbes immergรฉes ayant un unique point double (correspondant ร un point critique dโindice $`1`$ de $`F`$).
A chaque composante de type $`1`$ ou $`3`$ on associe un sommet de $`๐ข`$. Notons $`K`$ la rรฉunion des composantes de type $`1`$ ou $`3`$. Lโouvert $`SK`$ est une rรฉunion finie de cylindres diffรฉomorphes ร $`S^1\times `$. A chaque cylindre $`C`$ on associe une arรชte dont les extrรฉmitรฉs sont les sommets associรฉs aux composantes de niveaux de $`F`$ qui contiennent $`C`$. Nous avons une application naturelle $`p_๐ข:S๐ข`$, et, si $`H`$, on peut รฉcrire $`H=H_๐ขp_๐ข`$, oรน $`H_๐ข`$ est dรฉfinie sur $`๐ข`$.
Nous pouvons รฉlaguer le graphe $`๐ข`$ pour obtenir un graphe $`๐ข^{}`$, de la maniรจre suivante. Le graphe $`๐ข`$ possรจde des sommets de degrรฉ $`1`$ ou $`3`$. Si $`v`$ est un sommet de degrรฉ $`1`$ de $`๐ข`$, nous retirons $`v`$ ainsi que lโarรชte ร laquelle il รฉtait reliรฉ, pour obtenir un nouveau graphe. Ce faisant, nous pouvons crรฉer un sommet de degrรฉ $`2`$ (ou un autre sommet de degrรฉ $`1`$ ร partir de la seconde itรฉration de ce procรฉdรฉ). Rรฉpรฉtons ce procรฉdรฉ jusquโร obtenir un graphe $`๐ข^{}`$ qui ne possรจde plus que des points de degrรฉ $`2`$ ou $`3`$. Les sommets de degrรฉ $`3`$ de $`๐ข^{}`$ sont en nombre $`2g2`$. En effet la quantitรฉ
$$\underset{v}{}2degr\stackrel{ยด}{e}(v),$$
oรน la somme porte sur tous les sommets de $`๐ข`$, est รฉgale ร la caractรฉristique dโEuler de la surface, et reste constante au cours de lโรฉlagage. On note $`๐ฑ`$ lโensemble de ces $`2g2`$ sommets. Nous supposons dans le thรฉorรจme suivant que lโaire totale de la forme $`\omega `$ est รฉgale ร $`2g2`$.
###### Thรฉorรจme 2.
Si $`H`$ est dans $``$, nous avons :
$$๐๐ฉ_S(\phi _H^1)=_SH\omega \underset{v๐ฑ}{}H_๐ข(v).$$
Dans la derniรจre partie, nous proposons une construction รฉlรฉmentaire dโun quasi-morphisme homogรจne dรฉfini sur le revรชtement universel du groupe des diffรฉomorphismes hamiltoniens dโune variรฉtรฉ symplectique connexe fermรฉe $`(V,\omega )`$ pour laquelle $`[\omega ]=rc_1(V)`$. Ici, $`[\omega ]`$ dรฉsigne la classe de cohomologie de la forme symplectique et $`c_1(V)`$ la premiรจre classe de Chern de $`V`$. Pour les besoins de la cause, nous รฉcrivons notre hypothรจse sous la forme :
$$s[\omega ]=2c_1(V),$$
$`s`$ est un rรฉel non-nul. Nous calculons la restriction de ce quasi-morphisme sur les isotopies hamiltoniennes supportรฉes dans une boule ; son expression fait alors intervenir un quasi-morphisme $`\tau _{B,\omega }`$ introduit par J. Barge et ร. Ghys, dont la construction est rappelรฉe dans le texte. Si $`(f_t)`$ est une isotopie hamiltonienne dans $`V`$, on note $`\{f_t\}`$ lโรฉlรฉment du revรชtement universel du groupe $`\mathrm{Ham}(V,\omega )`$ quโelle dรฉfinit.
###### Thรฉorรจme 3.
Si $`(V,\omega )`$ vรฉrifie lโhypothรจse ci-dessus, il existe un quasi-morphisme homogรจne
$$๐:\stackrel{~}{\mathrm{Ham}}(V,\omega ),$$
tel que, pour toute isotopie hamiltonienne $`\{f_t\}`$ supportรฉe dans une boule $`B`$, on ait :
$$๐(\{f_t\})=\tau _{B,\omega }(f_1)+s๐๐ฉ_B(f_1).$$
Le texte est organisรฉ comme suit. Dans la seconde partie, nous prouvons les thรฉorรจmes $`1`$ et $`2`$, aprรจs avoir rappelรฉ des rรฉsultats de Banyaga. Dans la troisiรจme partie, nous rappelons des constructions de Barge et Ghys, puis nous prouvons le thรฉorรจme $`3`$. Exception faite des rรฉsultats mentionnรฉs dans le paragraphe $`2.1`$, qui sont utilisรฉs dans la troisiรจme partie, les parties $`2`$ et $`3`$ sont indรฉpendantes.
Une annonce des thรฉorรจmes $`1`$ et $`2`$ est contenue dans .
## 2. Surfaces de genre supรฉrieur
### 2.1. Extension du groupe des diffรฉomorphismes hamiltoniens
Les rรฉsultats de ce paragraphe sont dus ร Banyaga ; nous les rappelons succinctement.
Considรฉrons une variรฉtรฉ (fermรฉe, connexe) $`M`$ munie dโune forme de contact $`\alpha `$ dont le champ de Reeb $`X`$ est induit par une action libre du cercle $`/`$. La variรฉtรฉ $`V`$, quotient de $`M`$ par lโaction du cercle, porte une forme symplectique $`\omega `$ telle que $`\pi ^{}\omega =d\alpha `$, oรน $`\pi :MV`$ est la projection canonique.
Dans cette situation nous avons une extension centrale par $`/`$ du groupe des diffรฉomorphismes hamiltoniens de $`V`$. Dรฉcrivons dโabord cette extension au niveau des algรจbres de Lie de champs de vecteurs. Un รฉlรฉment $`Y`$ de lโalgรจbre $`_\alpha (M)`$ des champs de vecteurs sur $`M`$ qui prรฉservent $`\alpha `$ est invariant par lโaction du cercle $`/`$, il dรฉfinit un champ de vecteurs hamiltonien $`\pi _{}Y`$ sur $`V`$. Lโalgรจbre $`_\alpha (M)`$ est donc une extension centrale par $``$ de lโalgรจbre $`\mathrm{ham}(V,\omega )`$ des champs de vecteurs hamiltoniens sur $`V`$. Si $`Z\mathrm{ham}(V,\omega )`$ (avec $`\iota _Z\omega =dH_Z`$, $`_VH_Z\omega ^n=0`$, $`\mathrm{dim}V=2n`$), le champ de vecteurs $`\theta (Z)=\widehat{Z}(H_Z\pi )X`$ (oรน $`\widehat{Z}`$ est le relevรฉ horizontal de $`Z`$ : $`\alpha (\widehat{Z})=0`$) prรฉserve $`\alpha `$. Lโapplication $`Z\theta (Z)`$ est un morphisme dโalgรจbres de Lie qui scinde lโextension.
Notons $`G_\alpha (M)_0`$ le groupe des diffรฉomorphismes de $`M`$ qui prรฉservent $`\alpha `$, isotopes ร lโidentitรฉ via une isotopie qui prรฉserve $`\alpha `$. On a alors une extension centrale :
Si $`(f_t)`$ est une isotopie hamiltonienne on note $`\mathrm{\Theta }(f_t)`$ lโisotopie de $`M`$ obtenue en โintรฉgrantโ $`\theta `$. Sa classe dโhomotopie ne dรฉpend que de celle de $`(f_t)`$. Ainsi, si le groupe $`\mathrm{Ham}(V,\omega )`$ est simplement connexe, lโextension prรฉcรฉdente est scindรฉe. Grรขce au thรฉorรจme de Banyaga qui assure que le groupe $`\mathrm{Ham}(V,\omega )`$ est simple, la section qui scinde lโextension est unique. Si $`V`$ est une surface de genre supรฉrieur, le groupe $`\mathrm{Ham}(V,\omega )`$ est simplement connexe , et lโextension ci-dessus est canoniquement scindรฉe.
### 2.2. Construction du quasi-morphisme $`๐๐ฉ_S`$
On suppose donc que $`S`$ est une surface fermรฉe orientรฉe, de genre $`g`$ supรฉrieur ou รฉgal ร $`2`$, munie dโune forme symplectique dโaire totale $`2g2`$ (dโaprรจs un thรฉorรจme de Moser deux telles formes sont lโimage lโune de lโautre par un diffรฉomorphisme isotope ร lโidentitรฉ, le choix de $`\omega `$ est donc sans importance).
Nous noterons $`M`$ la variรฉtรฉ des droites orientรฉes tangentes ร $`S`$, $`\stackrel{~}{S}`$ le revรชtement universel de $`S`$, et $`\stackrel{~}{M}`$ la variรฉtรฉ des droites orientรฉes tangentes ร $`\stackrel{~}{S}`$. Le choix dโune mรฉtrique ร courbure constante sur $`S`$, de forme dโaire รฉgale ร $`\omega `$, fournit une structure de $`S^1`$-fibrรฉ principal sur $`M`$ et $`\stackrel{~}{M}`$. Nous noterons $`X`$ le champ de vecteurs sur $`M`$ tangent aux fibres de lโapplication $`\pi :MS`$ engendrรฉ par cette action du cercle. On note รฉgalement $`S_{\mathrm{}}^1`$ le cercle ร lโinfini de $`\stackrel{~}{S}`$ dรฉterminรฉ par cette mรฉtrique, et $`p_{\mathrm{}}:\stackrel{~}{M}S_{\mathrm{}}^1`$ la projection. La variรฉtรฉ $`\stackrel{~}{M}`$ est diffรฉomorphe ร $`\stackrel{~}{S}\times S_{\mathrm{}}^1`$ et le feuilletage $`(\stackrel{~}{S}\times \{\})_{S_{\mathrm{}}^1}`$ descend en un feuilletage sur $`M`$ appelรฉ feuilletage horocyclique. Nous noterons $`\alpha `$ une $`1`$-forme sur $`M`$ telle que $`\alpha (X)=1`$ et $`\pi ^{}\omega =d\alpha `$. La forme $`\alpha `$ est une forme de contact de champ de Reeb รฉgal ร $`X`$. On notera encore $`X`$ le relevรฉ ร $`\stackrel{~}{M}`$ de ce champ. Enfin, si $`\gamma :[\mathrm{0,1}]S_{\mathrm{}}^1`$ est un chemin continu, nous noterons $`n(\gamma )`$ lโentier dรฉfini comme suit. Si un paramรฉtrage de $`S_{\mathrm{}}^1`$ par $`/`$ est donnรฉ, notons $`\stackrel{~}{\gamma }`$ un relevรฉ de $`\gamma `$ ร $``$. On pose :
$$n(\gamma )=[\stackrel{~}{\gamma }(1)\stackrel{~}{\gamma }(0)],$$
cet entier ne dรฉpend pas du choix du paramรฉtrage. Si $`\gamma `$ et $`\beta `$ sont deux chemins dans $`S_{\mathrm{}}^1`$ avec $`\gamma (1)=\beta (0)`$ nous avons :
$$()|n(\gamma \beta )n(\gamma )n(\beta )|2.$$
Observations. Avant de construire le quasi-morphisme annoncรฉ dans le thรฉorรจme $`1`$, nous commenรงons par quelques remarques. A chaque ouvert connexe $`US`$, distinct de $`S`$, nous allons associer un รฉlรฉment canonique de lโespace
$$\frac{QM_h(\pi _1(U),)}{\mathrm{Hom}(\pi _1(U),)}.$$
Puisquโun quasi-morphisme homogรจne est invariant par conjugaison, nous oublierons parfois de choisir un point base pour le groupe $`\pi _1(U)`$. Soit $`\psi :U\times S^1\pi ^1(U)`$ une trivialisation du fibrรฉ $`\pi :MS`$ au-dessus de $`U`$. Si $`z_0`$ est un point sur $`S^1`$ et si $`\gamma `$ est un lacet dans $`U`$, basรฉ en $`x_0`$, on pose :
$$\varphi _{z_0}([\gamma ])=n(p_{\mathrm{}}(\stackrel{~}{\psi (\gamma (t),z_0)})),$$
$`\stackrel{~}{\psi (\gamma (t),z_0)}`$ est un relevรฉ ร $`\stackrel{~}{M}`$ du chemin $`\psi (\gamma (t),z_0)`$. Vu la propriรฉtรฉ $`()`$, lโapplication $`\varphi _{z_0}`$ est un quasi-morphisme sur le groupe $`\pi _1(U,x_0)`$. Son homogรฉnรฉisรฉ $`\varphi `$ ne dรฉpend pas du point $`z_0`$. Si deux mรฉtriques ร courbure constante sur $`S`$ sont donnรฉes, on peut identifier par un homรฉomorphisme les bords ร lโinfini de $`\stackrel{~}{S}`$ associรฉs, lโindice $`n`$ ne dรฉpend donc pas de la mรฉtrique. A lโaddition dโun homomorphisme prรจs, il ne dรฉpend pas de la trivialisation $`\psi `$, la classe
$$[\varphi ]\frac{QM_h(\pi _1(U),)}{\mathrm{Hom}(\pi _1(U),)}$$
est donc canonique. Une autre maniรจre de la dรฉcrire serait la suivante. On fixe un point $`\stackrel{~}{x_0}M`$ au-dessus de $`x_0`$. Si $`[\gamma ]\pi _1(U,x_0)`$, on note $`\stackrel{~}{\gamma }`$ son relevรฉ issu de $`\stackrel{~}{x_0}`$ tangent au feuilletage horocyclique. On peut รฉcrire $`\stackrel{~}{\gamma }(t)=\psi (\gamma (t),z(t))`$. Notons $`f([\gamma ])`$ la variation de lโargument de $`z(t)`$ comptรฉe en tours. Alors $`f`$ est un quasi-morphisme dont lโhomogรฉnรฉisรฉ est รฉgal ร $`\varphi `$. On peut en quelque sorte penser ร un feuilletage transverse aux fibres du fibrรฉ $`MS`$ comme ร une โquasi-trivialisationโ du fibrรฉ. Alors que la comparaison de deux trivialisations du fibrรฉ au-dessus de lโouvert $`U`$ fournit un homomorphisme $`\pi _1(U)`$, la comparaison du feuilletage avec une trivialisation fournit un quasi-morphisme homogรจne sur le groupe $`\pi _1(U)`$. Si le groupe $`\pi _1(U)`$ est abรฉlien, il nโadmet pas de quasi-morphisme homogรจne non-trivial (cโest-ร -dire autre que les homomorphismes). En revanche si le groupe $`\pi _1(U)`$ est libre non-abรฉlien, la classe $`[\varphi ]`$ va apparaรฎtre comme une obstruction ร ce que le quasi-morphisme $`๐๐ฉ_S:\mathrm{Ham}(S,\omega )`$ que nous allons dรฉfinir se restreigne en le morphisme de Calabi sur le groupe $`\mathrm{\Gamma }_U`$.
Donnons enfin une derniรจre interprรฉtation de la classe $`[\varphi ]`$. Rappelons quโร toute reprรฉsentation $`\rho `$ dโun groupe discret $`\mathrm{\Gamma }`$ dans le groupe $`\mathrm{Hom}\stackrel{ยด}{\mathrm{e}}\mathrm{o}_+(S^1)`$ des homรฉomorphismes du cercle qui prรฉservent lโorientation, on peut associer une classe de cohomologie bornรฉe $`eu_b(\rho )H_b^2(\mathrm{\Gamma },)`$, appelรฉe classe dโEuler bornรฉe, voir . La reprรฉsentation $`\pi _1(S)\mathrm{PSL}(2,)`$ associรฉe ร une mรฉtrique ร courbure constante sur $`S`$ fournit donc une classe $`eu_b(S)H_b^2(\pi _1(S),)`$, qui ne dรฉpend pas de la mรฉtrique. Pour tout groupe discret $`\mathrm{\Gamma }`$, le noyau de lโapplication $`H_b^2(\mathrm{\Gamma },)H^2(\mathrm{\Gamma },)`$ est isomorphe ร lโespace
$$\frac{QM_h(\mathrm{\Gamma },)}{\mathrm{Hom}(\mathrm{\Gamma },)}.$$
Si $`US`$ est un ouvert connexe distinct de $`S`$, on note $`i_U:\pi _1(U)\pi _1(S)`$ le morphisme naturel. Puisque le groupe $`H^2(\pi _1(U),)`$ est trivial, la classe $`i_U^{}eu_b(S)`$ (que nous considรฉrons comme une classe rรฉelle) est dans le noyau
$$\mathrm{Ker}(H_b^2(\pi _1(U),)H^2(\pi _1(U),)).$$
Cโest la classe $`[\varphi ]`$ prรฉcรฉdemment dรฉcrite.
Nous pouvons maintenant construire le quasi-morphisme $`๐๐ฉ_S`$. Soit $`(f_t)`$ une isotopie hamiltonienne dans $`S`$. Notons $`\mathrm{\Theta }(f_t):MM`$ lโisotopie qui relรจve $`(f_t)`$ prรฉcรฉdemment construite, et $`(F_t)`$ lโisotopie de $`\stackrel{~}{M}`$ qui relรจve $`\mathrm{\Theta }(f_t)`$. Si $`v`$ et $`w`$ sont deux points de $`\stackrel{~}{M}`$ tels que $`\stackrel{~}{\pi }(v)=\stackrel{~}{\pi }(w)`$ (oรน $`\stackrel{~}{\pi }`$ est la projection $`\stackrel{~}{M}\stackrel{~}{S}`$), on รฉcrit $`w=\varphi _X^{u_0}(v)`$ avec $`u_0[\mathrm{0,1}]`$, oรน $`\varphi _X^u`$ est le flot de $`X`$. Notons $`G(u,t)=p_{\mathrm{}}(F_t(\varphi _X^{uu_0}v))`$. Le lacet lu sur le โbordโ de $`G`$ a un indice $`n`$ รฉgal ร $`0`$. Puisque les chemins $`G(\mathrm{,0})`$ et $`G(\mathrm{,1})`$ ont un indice bornรฉ par $`1`$, on a $`|n(p_{\mathrm{}}(F_t(v)))n(p_{\mathrm{}}(F_t(w))|2`$. On dรฉfinit alors, pour $`\stackrel{~}{x}\stackrel{~}{S}`$, $`\stackrel{~}{\mathrm{angle}}(\stackrel{~}{x},f_1)=\mathrm{inf}_{\stackrel{~}{\pi }(v)=\stackrel{~}{x}}n(p_{\mathrm{}}(F_t(v)))`$. La fonction $`\stackrel{~}{\mathrm{angle}}(,f_1)`$ est invariante sous lโaction du groupe fondamental de $`S`$ et dรฉfinit une fonction mesurable bornรฉe $`\mathrm{angle}(,f_1)`$ sur $`S`$. Elle vรฉrifie :
$$|\mathrm{angle}(x,f_1g_1)\mathrm{angle}(x,g_1)\mathrm{angle}(g_1(x),f_1)|8.$$
Ainsi lโapplication qui, au diffรฉomorphisme hamiltonien $`f_1`$ associe lโintรฉgrale
$$_S\mathrm{angle}(,f_1)\omega $$
est un quasi-morphisme. Nous pouvons lโhomogรฉnรฉiser pour dรฉfinir :
$$๐๐ฉ_S(f_1)=\mathrm{lim}_p\mathrm{}\frac{1}{p}_S\mathrm{angle}(,f_1^p)\omega .$$
Dโaprรจs le thรฉorรจme ergodique sous-additif , la suite de fonctions
$$\frac{1}{p}\mathrm{angle}(,f^p)$$
converge $`\omega `$-presque partout quand $`p`$ tend vers lโinfini, vers une fonction mesurable $`\widehat{\mathrm{angle}}(,f)`$. Il nโest pas difficile de vรฉrifier que cette fonction est bornรฉe et satisfait :
$$๐๐ฉ_S(f)=_S\widehat{\mathrm{angle}}(,f)\omega .$$
Discutons maintenant des diffรฉrents choix effectuรฉs pour notre construction. Nous avons dรฉjร indiquรฉ pourquoi lโindice $`n`$ dโune courbe ne dรฉpend pas du choix de la mรฉtrique. Par ailleurs :
* Si lโon change de mรฉtrique, lโaction du cercle sur $`M`$ (i.e. le champ $`X`$) change, mais la classe dโEuler du fibrรฉ nโรฉtant pas modifiรฉe, on peut trouver un diffรฉomorphisme de $`M`$, induisant lโidentitรฉ sur $`S`$, qui entrelace les deux actions. On en dรฉduit aisรฉment lโinvariance du quasi-morphisme.
* Lorsque la mรฉtrique est fixรฉe, le choix de la forme $`\alpha `$ est sans importance. Une autre primitive de $`\pi ^{}\omega `$ valant $`1`$ sur $`X`$ serait de la forme $`\alpha +\pi ^{}\beta `$, oรน $`\beta `$ est une $`1`$-forme fermรฉe sur la surface. En utilisant la nullitรฉ du flux dโune isotopie hamiltonienne, on voit que le quasi-morphisme final est inchangรฉ. Notons que, lorsque $`M`$ est identifiรฉ au fibrรฉ unitaire tangent ร $`S`$, la forme $`\alpha `$ peut-รชtre prise nulle dans la direction du flot gรฉodรฉsique.
* Une fois acquise lโindรฉpendance de $`๐๐ฉ_S`$ vis-ร -vis de la mรฉtrique, lโinvariance par conjugaison dans $`\mathrm{Symp}(S,\omega )`$ est claire. Il suffit de considรฉrer une mรฉtrique (ร courbure constante, de forme dโaire $`\omega `$) et de la transporter par le diffรฉomorphisme symplectique considรฉrรฉ.
Supposons que $`US`$ soit un ouvert connexe distinct de $`S`$ et $`(f_t)`$ une isotopie hamiltonienne dans $`U`$. Nous allons calculer $`๐๐ฉ_S(f_1)`$.
Choisissons une trivialisation $`\psi :U\times S^1\pi ^1(U)`$ du fibrรฉ $`M`$ au-dessus de $`U`$, telle que $`X=\frac{}{s}`$ (oรน $`s`$ est la coordonnรฉe angulaire sur le cercle). On a alors $`\alpha =ds+\pi ^{}\lambda `$, oรน $`\lambda `$ est une primitive de $`\omega `$ sur $`U`$. Notons รฉgalement $`\varphi `$ le quasi-morphisme homogรจne sur $`\pi _1(U)`$ associรฉ ร cette trivialisation. Enfin, on note $`Z_t`$ le champ de vecteurs qui engendre lโisotopie, $`H_t`$ un hamiltonien pour $`Z_t`$ avec $`\mathrm{supp}(H_t)U`$ et $`\stackrel{~}{H}_t`$ la fonction dโintรฉgrale nulle sur $`S`$ qui diffรจre de $`H_t`$ par une constante.
Fixons un point $`x_0`$ dans $`U`$ et un compact $`K`$ tel que $`\mathrm{supp}(f_t)K`$. Pour tout point $`x`$ de $`K`$, on choisit un chemin $`\alpha _{x_0x}`$ de classe $`C^1`$ par morceaux de $`x_0`$ ร $`x`$, contenu dans $`U`$, de dรฉrivรฉe bornรฉe indรฉpendamment de $`x`$. On note $`\gamma _{x,f}`$ le lacet $`\alpha _{x_0x}(f_t(x))\overline{\alpha _{x_0f(x)}}`$ et $`[\varphi ],f(x)`$ la limite de la suite $`(\frac{1}{p}\varphi ([\gamma _{x,f^p}]))_{p0}`$ (pour relier $`f^p`$ ร lโidentitรฉ on utilise bien sรปr lโisotopie $`(f_t)`$ concatรฉnรฉe $`p`$ fois).
###### Proposition 2.1.
Pour presque tout $`x`$ de $`U`$ nous avons :
$$\widehat{\mathrm{angle}}(x,f)=[\varphi ],f(x)+๐(y_0^1(\lambda (Z_t)+\stackrel{~}{H}_t)(f_t(y))๐t)(x).$$
Dans cette proposition, et dans la suite, $`๐(\phi )`$ dรฉsignera la limite des moyennes de Birkhoff dโune fonction (intรฉgrable) $`\phi `$ relativement ร la transformation $`f=f_1`$.
Preuve : Notons $`(h_t)`$ lโisotopie obtenue en concatรฉnant $`p`$ fois $`(f_t)`$. Si $`x`$ est dans $`U`$ et $`z`$ dans $`S^1`$ nous avons :
$$\mathrm{\Theta }(h_t)(\psi (x,z))=\psi (h_t(x),exp(2i\pi _0^t(\lambda (Z_t^{})+\stackrel{~}{H}_t^{})(h_t^{}(x))๐t^{})z).$$
Notons $`v(t)`$, $`v_1(t)`$, $`v_2(t)`$ les courbes
$$\mathrm{\Theta }(h_t)(\psi (x,z)),$$
$$\psi (h_t(x),z),$$
$$\psi (f^p(x),exp(2i\pi _0^t(\lambda (Z_t^{})+\stackrel{~}{H}_t^{})(h_t^{}(x))๐t^{})),$$
respectivement. On choisit deux relevรฉs $`\stackrel{~}{v}_1`$ et $`\stackrel{~}{v}_2`$ ร $`\stackrel{~}{M}`$ tels que $`\stackrel{~}{v}_1(1)=\stackrel{~}{v}_2(0)`$, et lโon note $`\stackrel{~}{v}=\stackrel{~}{v}_1\stackrel{~}{v}_2`$. Dans la suite dโรฉgalitรฉs ci-dessous, le symbole $``$ voudra dire que les deux membres de lโรฉgalitรฉ diffรจrent dโune quantitรฉ bornรฉe, la valeur exacte de la borne important peu (mais pouvant cependant aisรฉment รชtre dรฉterminรฉe).
$$\mathrm{angle}(x,f^p)n(p_{\mathrm{}}(\stackrel{~}{v}(t)))n(p_{\mathrm{}}(\stackrel{~}{v}_1(t)))n(p_{\mathrm{}}(\stackrel{~}{v}_2(t))),$$
$$n(p_{\mathrm{}}(\stackrel{~}{v}_2(t)))_0^1(\lambda (Z_t^{})+\stackrel{~}{H}_t^{})(h_t^{}(x))๐t^{},$$
$$n(p_{\mathrm{}}(\stackrel{~}{v}_1(t)))\varphi ([\gamma _{x,f^p}]).$$
Nous obtenons au total lโexistence dโune constante $`C`$ telle que :
$$|\mathrm{angle}(x,f^p)_0^1(\lambda (Z_t^{})+\stackrel{~}{H}_t^{})(h_t^{}(x))๐t^{}+\varphi ([\gamma _{x,f^p}])|C.$$
Le rรฉsultat suit. $`\mathrm{}`$
Une fois la proposition prรฉcรฉdente acquise, nous pouvons conclure la preuve du thรฉorรจme $`1`$. Hors de $`U`$, la fonction $`\widehat{\mathrm{angle}}(x,f)`$ est รฉgale ร $`_0^1\stackrel{~}{H}_t(x)๐t`$. Nous obtenons donc :
$$๐๐ฉ_S(f)=_U[\varphi ],f\omega +๐๐ฉ_U(f)$$
(la fonction $`[\varphi ],f`$ est bien sรปr nulle hors de lโouvert $`U`$). Si $`U`$ est simplement connexe, le terme $`[\varphi ],f`$ est identiquement nul. Si $`U`$ est un anneau, le quasi-morphisme homogรจne $`\varphi `$ est un morphisme qui se reprรฉsente par une $`1`$-forme fermรฉe $`\beta `$. Dans ce cas, on vรฉrifie sans peine que
$$[\varphi ],f(x)=๐(y_0^1\beta (Z_t)(f_t(y))๐t)(x).$$
On en dรฉduit :
$$_U[\varphi ],f\omega =_U_0^1\beta (Z_t)๐t\omega .$$
Mais cette derniรจre intรฉgrale est nulle pour une isotopie hamiltonienne. Dans le cas dโun flot hamiltonien, cela traduit simplement le fait que le cycle asymptotique de Schwartzman pour la mesure $`\omega `$ est nul. On a donc bien $`๐๐ฉ_S(f_1)=๐๐ฉ_U(f_1)`$ dans ce cas. Si le groupe $`\pi _1(U)`$ est libre non-abรฉlien, lโespace
$$\frac{QM_h(\pi _1(U),)}{\mathrm{Hom}(\pi _1(U),)}$$
nโest pas trivial , et lโintรฉgrale $`_U[\varphi ,f]\omega `$ peut ne pas sโannuler. Nous avons achevรฉ la preuve du thรฉorรจme $`1`$.
Dรฉcrivons maintenant les constructions de Gambaudo et Ghys qui permettent de prouver que lโespace des quasi-morphismes homogรจnes sur $`\mathrm{Ham}(S,\omega )`$ nuls en restriction aux sous-groupes $`(\mathrm{\Gamma }_U)`$$`U`$ est un disque ou un anneau, est de dimension infinie. Considรฉrons une $`1`$-forme $`\eta `$, non nรฉcessairement fermรฉe, sur la surface $`S`$, et notons $`\stackrel{~}{\eta }`$ son relรจvement ร $`\stackrel{~}{S}`$. On suppose toujours quโune mรฉtrique ร courbure constante de forme dโaire $`\omega `$ est fixรฉe. Nous allons construire un quasi-morphisme homogรจne
$$\mathrm{\Phi }_\eta :\mathrm{Symp}_0(S,\omega ).$$
Un รฉlรฉment $`f`$ de $`\mathrm{Symp}_0(S,\omega )`$ admet un relevรฉ canonique $`\stackrel{~}{f}`$ ร $`\stackrel{~}{S}`$ : puisque le groupe $`\mathrm{Symp}_0(S,\omega )`$ est simplement connexe nous choisissons une isotopie $`(f_t)`$ reliant lโidentitรฉ ร $`f`$, et nous considรฉrons son relevรฉ $`(\stackrel{~}{f}_t)`$ ร $`\stackrel{~}{S}`$. On pose alors $`\stackrel{~}{f}=\stackrel{~}{f}_1`$. Si $`x`$ est dans $`\stackrel{~}{S}`$ nous noterons $`\delta (x,f)`$ lโunique gรฉodรฉsique (pour la mรฉtrique fixรฉe) qui relie $`x`$ ร $`\stackrel{~}{f}(x)`$. La fonction
$$x_{\delta (x,f)}\stackrel{~}{\eta },$$
est $`\pi _1(S)`$-invariante et dรฉfinit donc une fonction $`u(\eta ,f)`$ sur la surface $`S`$. En utilisant le fait que la $`2`$-forme $`d\stackrel{~}{\eta }`$ vรฉrifie une inรฉgalitรฉ de la forme $`d\stackrel{~}{\eta }C\stackrel{~}{\omega }`$ sur $`\stackrel{~}{S}`$, et le fait que les triangles gรฉodรฉsiques de $`\stackrel{~}{S}`$ sont dโaire bornรฉe, nous obtenons aisรฉment :
$$|u(\eta ,fg)u(\eta ,g)u(\eta ,f)g|\pi C.$$
Ainsi $`f_Su(\eta ,f)\omega `$ dรฉfinit un quasi-morphisme sur $`\mathrm{Symp}_0(S,\omega )`$, son homogรฉnรฉisรฉ sera notรฉ $`\mathrm{\Phi }_\eta `$ (contrairement ร $`๐๐ฉ_S`$, il dรฉpend de la mรฉtrique).
Dans , il est montrรฉ que la famille $`(\mathrm{\Phi }_\eta )_\eta `$, restreinte au groupe $`\mathrm{Ham}(S,\omega )`$, engendre un espace de dimension infinie dans lโespace $`QM_h(\mathrm{Ham}(S,\omega ),)`$. Les auteurs utilisent pour cela des โtwistsโ supportรฉs au voisinage dโune gรฉodรฉsique fermรฉe simple. Notons que, pour que ceux-ci soient hamiltoniens, cette gรฉodรฉsique doit รชtre homologue ร $`0`$.
Supposons que $`(f_t)`$ est une isotopie symplectique supportรฉe dans le compact $`K`$ contenu dans lโouvert simplement connexe $`U`$ de $`S`$. Notons $`\stackrel{~}{K}\stackrel{~}{U}`$ des relevรฉs ร $`\stackrel{~}{S}`$. Le compact $`\stackrel{~}{K}`$ est stable par $`\stackrel{~}{f}`$ et de diamรจtre fini. On a donc $`|u(\eta ,f^p)||\eta |\mathrm{diam}(\stackrel{~}{K})`$ pour tout $`p`$. Ceci assure que $`\mathrm{\Phi }_\eta (f_1)=0`$. Soit maintenant $`(f_t)`$ une isotopie hamiltonienne (engendrรฉe par le champ $`Z_t`$) supportรฉe dans lโanneau $`A=]\mathrm{0,1}[\times /S`$. On suppose bien sรปr ce plongement injectif au niveau du groupe fondamental, sans quoi on serait ramenรฉ au cas prรฉcรฉdent. On vรฉrifie alors que
$$\widehat{u}(\eta ,f)=\mathrm{lim}_p\mathrm{}\frac{1}{p}u(\eta ,f^p)$$
vaut $`l๐(_0^1\beta (Z_t)f_t๐t)`$, oรน la classe de la $`1`$-forme fermรฉe $`\beta `$ sur $`A`$ engendre $`H^1(A,)`$ et $`l`$ est lโintรฉgrale de la $`1`$-forme $`\eta `$ sur la gรฉodรฉsique fermรฉe librement homotope dans $`S`$ au gรฉnรฉrateur de $`\pi _1(A)`$. On en dรฉduit $`\mathrm{\Phi }_\eta (f_1)=0`$.
### 2.3. Calcul sur des hamiltoniens autonomes
Nous prouvons ici le thรฉorรจme $`2`$. Notons $`U`$ lโouvert $`S\{x_l\}`$ et fixons une trivialisation du fibrรฉ $`\pi :MS`$ au-dessus de $`U`$. Celle-ci fournit une primitive $`\lambda `$ de $`\omega `$ sur $`U`$ et un quasi-morphisme homogรจne $`\varphi `$ sur le groupe $`\pi _1(U)`$, comme prรฉcรฉdemment. Si $`a`$ est une arรชte du graphe $`๐ข`$ nous noterons $`a^+`$ et $`a^{}`$ les sommets ร ses extrรฉmitรฉs (avec la convention $`F_๐ข(a^{})<F_๐ข(a^+)`$). Pour chaque arรชte $`a`$, nous fixons un paramรฉtrage de $`p_๐ข^1(a)`$ par $`/\times ]F_๐ข(a^{}),F_๐ข(a^+)[`$, avec des coordonnรฉes $`(\theta ,t)`$ telles que $`F(\theta ,t)=t`$ et $`\omega =d\theta dt`$. Si $`H`$ est une fonction dans $``$, le champ hamiltonien $`Z_H`$ sโรฉcrit sur $`p_๐ข^1(a)`$ :
$$\vartheta (t)\frac{}{\theta },$$
oรน la fonction $`\vartheta `$ satisfait :
$$_{F_๐ข(a^{})}^{F_๐ข(a^+)}\vartheta (t)๐t=H_๐ข(a^+)H_๐ข(a^{}).$$
Pour un domaine ร bord lisse $`DU`$ nous noterons $`[\varphi ],D`$ la somme des valeurs de $`\varphi `$ sur les classes de conjugaison dรฉterminรฉes par chacune des composantes du bord de $`D`$. Cette valeur ne dรฉpend que de la classe $`[\varphi ]`$.
###### Proposition 2.2.
Si $`D`$ est ร bord gรฉodรฉsique, pour une mรฉtrique ร courbure constante quelconque sur $`S`$, on a $`[\varphi ],D=\chi (D)`$.
Preuve : on peut supposer que $`\omega `$ est la forme dโaire associรฉe ร la mรฉtrique qui rend le bord de $`D`$ gรฉodรฉsique (car le quasi-morphisme $`\varphi `$ est indรฉpendant de $`\omega `$). On a vu que la $`1`$-forme $`\alpha `$ peut alors รชtre choisie nulle dans la direction du flot gรฉodรฉsique. Au-dessus de $`U`$, dans une trivialisation dans laquelle $`X=/s`$, on a $`\alpha =ds+\pi ^{}(\lambda )`$. Soit $`\gamma `$ une orbite pรฉriodique du flot gรฉodรฉsique telle que $`\pi (\gamma )`$ est une composante du bord de $`D`$. Puisque $`\gamma `$ est fermรฉe on a : $`\varphi ([\pi (\gamma )])=_\gamma ๐s=_{\pi (\gamma )}\lambda `$. En sommant sur les diffรฉrentes composantes de bord, on obtient : $`[\varphi ],D=_D\omega =\chi (D)`$. $`\mathrm{}`$
Si $`H`$ est dans $``$, bien que lโisotopie $`(\phi _H^t)`$ ne soit pas nรฉcessairement ร support dans $`U`$, on peut rรฉpรฉter le raisonnement qui a servi ร รฉtablir la proposition $`2.1`$. Si $`x`$ est dans $`U`$ et $`p_๐ข(x)`$ nโest pas un sommet de $`๐ข`$, on notera $`[x]`$ la classe dโhomotopie libre du cercle $`p_๐ข^1(p_๐ข(x))`$, orientรฉ par $`X_F`$. Notant $`\stackrel{~}{H}`$ la fonction dโintรฉgrale nulle sur $`S`$ qui diffรจre de $`H`$ par une constante, on a, presque partout sur $`U`$ :
$$\widehat{\mathrm{angle}}(x,\phi _H^1)=๐(y\lambda (X_H)(y)+\stackrel{~}{H}(y))(x)\vartheta (x)\varphi ([x]).$$
La fonction $`๐(y\lambda (X_H)(y)+\stackrel{~}{H}(y))`$ a pour intรฉgrale $`_SH\omega (2g2)H(x_l)`$. Si $`a`$ est une arรชte de $`๐ข`$, nous noterons $`[a]`$ la valeur commune des classes $`[x]`$ pour $`xp_๐ข^1(a)`$. La somme des intรฉgrales de $`\vartheta (x)\varphi ([x])`$ sur les diffรฉrentes arรชtes vaut :
$$\underset{a}{}\varphi ([a])(H_๐ข(a^+)H_๐ข(a^{})).$$
On peut รฉcrire la somme sous la forme
$$\underset{v}{}C(v)H_๐ข(v),$$
oรน la somme porte cette fois sur les sommets de $`๐ข`$. Il faut calculer les constantes $`C(v)`$. Au prรฉalable, notons le fait suivant :
Observation. Si $`x`$ et $`y`$ sont deux lacets dans $`\pi ^1(U)`$, tangents au feuilletage horocyclique de $`S`$, ayant mรชme point base, alors $`\varphi ([\pi (xy)])=\varphi ([\pi (x)])+\varphi ([\pi (y)])`$.
On peut alors commencer le calcul des constantes $`C(v)`$ :
* Si $`v`$ correspond ร un extremum local autre que le maximum global, la constante $`C(v)`$ est nulle. En effet, elle est รฉgale ร la valeur de $`\varphi `$ sur un petit lacet qui entoure lโextremum. Comme celui-ci est dans $`U`$ le lacet est nul dans $`\pi _1(U)`$.
* Si $`v=p_๐ข(x_l)`$, la constante $`C(v)`$ est รฉgale ร la valeur de $`\varphi `$ sur la classe dโhomotopie dans $`U`$ dโun petit lacet qui entoure $`x_l`$ (avec lโorientation opposรฉe ร celle du bord de $`\{F\lambda _lฯต\}`$). Si ce lacet $`\gamma `$ est assez petit, il admet un relevรฉ tangent au feuilletage horocyclique qui est fermรฉ. On constate alors que $`\varphi ([\gamma ])=(2g2)`$ (la classe dโEuler du fibrรฉ en cercles au-dessus de $`S`$).
Il reste alors ร calculer les constantes $`C(v)`$ pour les sommets $`v`$ correspondant ร des points critiques dโindice $`1`$. Il nโest pas difficile de vรฉrifier que
$$C(p_๐ข(x_j))=[\varphi ],\{\lambda _jฯตF\lambda _j+ฯต\},$$
(pour $`ฯต`$ assez petit). Le domaine $`\{\lambda _jฯตF\lambda _j+ฯต\}`$ est constituรฉ dโun nombre fini de cylindres sur le bord desquels $`\varphi `$ est nul et dโun pantalon $`P`$. Il faut รฉvaluer le terme $`[\varphi ],P`$.
* Si lโune des composantes du bord de $`P`$ est homotope ร $`0`$ dans $`U`$, $`\varphi `$ est nul รฉvaluรฉ contre celle-ci. Les deux autres composantes sont alors librement homotopes dans $`U`$, et les deux valeurs de $`\varphi `$ correspondantes sont opposรฉes. On peut donc supposer les trois composantes de $`P`$ essentielles dans $`U`$ sans quoi $`C(v)=0`$.
* Si en outre ces trois composantes sont essentielles dans $`S`$, on peut trouver une mรฉtrique ร courbure constante (de forme dโaire รฉgale ร $`\omega `$) qui rend le bord de $`P`$ gรฉodรฉsique (ou seulement deux composantes sur trois de $`P`$ si deux dโentre elles sont librement homotopes). Une gรฉnรฉralisation immรฉdiate de la proposition $`2.2`$ permet alors de montrer que la constante $`C(v)`$ vaut $`1`$.
* Il reste ร traiter le cas oรน lโune des composantes, disons $`\alpha _1`$, de $`P`$ est contractile dans $`S`$. Dans ce cas elle borde un disque plongรฉ qui contient le point $`x_l`$ (puisque lโon a supposรฉ cette mรชme courbe essentielle dans $`U`$). Les deux autres composantes $`\alpha _2`$ et $`\alpha _3`$ de $`P`$ sont alors essentielles dans $`S`$. Modifions le pantalon $`P`$ en un pantalon $`P^{}`$ de composantes de bord $`(\alpha _i^{})_{1i3}`$ telles que $`\alpha _i^{}`$ est homotope ร $`\alpha _i`$, et $`\alpha _1^{}`$ et $`\alpha _2^{}`$ ont mรชme point base. Le lacet $`\alpha _1^{}`$ peut รชtre choisi contenu dans un disque arbitrairement petit au voisinage de $`x_l`$. On peut alors trouver une mรฉtrique ร courbure constante qui rend $`\alpha _2^{}`$ gรฉodรฉsique. Notant $`\beta _2`$ lโorbite pรฉriodique du flot gรฉodรฉsique telle que $`\pi (\beta _2)=\alpha _2^{}`$, on peut trouver un relevรฉ, tangent au feuilletage horocyclique et fermรฉ, $`\beta _1`$ de $`\alpha _1^{}`$, issu du mรชme point que lโorbite $`\beta _2`$. Dans ce cas lโobservation ci-dessus assure que $`\varphi ([\alpha _1^{}\alpha _2^{}])=\varphi ([\alpha _1^{}])+\varphi ([\alpha _2^{}])`$. Puisque la derniรจre composante de bord de $`P^{}`$ dรฉfinit la classe de conjugaison de $`[\alpha _1^{}\alpha _2^{}]^1`$, on a $`[\varphi ],P^{}=0`$.
Il nโest pas difficile de vรฉrifier que lโensemble des sommets pour lesquels $`C(v)=1`$ correspond ร lโensemble $`๐ฑ`$ dรฉfini dans lโintroduction. Une autre maniรจre de dรฉcrire cet ensemble serait la suivante. Les รฉlรฉments de $`๐ฑ`$ sont les sommets correspondant ร des points critiques dโindice $`1`$ pour lesquels les trois composantes du bord du pantalon $`P`$ dรฉcrit prรฉcรฉdemment sont essentielles dans $`S`$. En rรฉsumรฉ, la constante $`C(p_๐ข(x_l))`$ vaut $`(2g2)`$, les autres constantes sont รฉgales ร $`1`$ pour les รฉlรฉments de $`๐ฑ`$ et $`0`$ sinon. Finalement, la somme initiale est รฉgale ร :
$$\underset{v๐ฑ}{}H_๐ข(v)(2g2)H(x_l).$$
Nous obtenons ainsi :
$$๐๐ฉ_S(\phi _H^1)=_S\widehat{\mathrm{angle}}(,\phi _H^1)\omega =_SH\omega \underset{v๐ฑ}{}H_๐ข(v).$$
## 3. Quasi-morphisme sur certaines variรฉtรฉs symplectiques de premiรจre classe de Chern non-nulle
### 3.1. Exposant de Lyapunov symplectique
Notons $`B`$ une boule de $`^{2n}`$, munie dโune forme symplectique $`\nu `$ de volume fini. Nous rappelons ici une construction due ร Barge et Ghys , dโun quasi-morphisme homogรจne $`\tau _{B,\nu }`$ sur le groupe $`\mathrm{\Gamma }_{B,\nu }=\mathrm{Diff}^c(B,\nu )`$ des diffรฉomorphismes symplectiques de $`B`$ ร support compact. Cette construction apparaรฎt dรฉjร implicitement dans . Nous renvoyons ร pour plus de dรฉtails.
On commence par construire un quasi-morphisme $`\mathrm{\Phi }`$ sur le revรชtement universel $`\stackrel{~}{\mathrm{Sp}(\mathrm{E},\omega )}`$ du groupe symplectique dโun espace vectoriel symplectique $`(E,\omega )`$. Cette construction est bien connue ; nous la rappelons cependant, dans un souci de complรฉtude. Supposons que $`J`$ soit une structure presque-complexe sur $`E`$ compatible avec $`\omega `$ : $`J`$ est un endomorphisme de $`E`$, de carrรฉ $`1`$, qui prรฉserve $`\omega `$, avec $`\omega (u,Ju)>0`$ pour tout vecteur non-nul $`u`$ de $`E`$. Muni de la forme $`(u,v)_J=\omega (u,Jv)i\omega (u,v)`$, $`E`$ devient un espace vectoriel hermitien. Nous noterons $`\mathrm{\Lambda }(E)`$ la grassmannienne lagrangienne de $`E`$. Si $`L_0`$ et $`L_1`$ sont deux lagrangiens, il existe un endomorphisme unitaire $`u`$ de $`E`$ tel que $`u(L_0)=L_1`$. Le nombre complexe
$$det_{}^2(u)S^1,$$
ne dรฉpend pas du choix de $`u`$. On le note $`det_{L_0}^2L_1`$. Il vรฉrifie la relation (de cocycle) :
$$det_{L_0}^2L_2=det_{L_0}^2L_1det_{L_1}^2L_2.$$
En particulier, si $`(L_t)`$ est une courbe dans $`\mathrm{\Lambda }(E)`$ la variation de lโargument du nombre complexe $`det_W^2L_t`$ (comptรฉe en tours) ne dรฉpend pas du choix de $`W`$. On note $`\mathrm{\Delta }(det^2L_t)`$ ce nombre.
###### Proposition 3.1.
Si $`(L_t)`$ est une courbe dans $`\mathrm{\Lambda }(E)`$ qui reste toujours transverse ร un lagrangien donnรฉ $`W`$, on a $`|\mathrm{\Delta }(det^2L_t)|n`$.
Il est classique que lโapplication
$$\begin{array}{ccc}\hfill \mathrm{\Lambda }(E)& & S^1\hfill \\ \hfill L& & det_{L_0}^2L\hfill \end{array}$$
induit un isomorphisme entre les groupes fondamentaux. Par ailleurs, il est รฉgalement bien connu que, pour tout lagrangien $`W`$, lโintersection avec lโhypersurface
$$\{L,\mathrm{dim}LW>0\},$$
engendre le groupe $`H^1(\mathrm{\Lambda }(E),)`$. Il nโest donc pas surprenant que le fait de rester transverse ร un lagrangien donnรฉ, empรชche une courbe de $`\mathrm{\Lambda }(E)`$ deโtrop tournerโ (voir ). Prouvons maintenant la proposition.
Preuve : Lโapplication qui a un endomorphisme $`f`$ de $`W`$, symรฉtrique pour le produit scalaire $`u,v_J=\omega (u,Jv)`$, associe le graphe $`L_f=\{f(x)+Jx\}_{xW}`$ est un diffรฉomorphisme sur lโouvert des lagrangiens transverses ร $`W`$. On a :
$$det_W^2L_f=\underset{k=1}{\overset{n}{}}\frac{(\lambda _k+i)^2}{1+\lambda _k^2},$$
oรน les $`\lambda _k`$ sont les valeurs propres de $`f`$. Si $`(f_t)`$ est un chemin dโendomorphismes symรฉtriques de $`W`$, de valeurs propres $`\lambda _1(t)\mathrm{}\lambda _n(t)`$, la variation de lโargument du nombre complexe
$$\underset{k=1}{\overset{n}{}}\frac{(\lambda _k(t)+i)^2}{1+\lambda _k(t)^2}$$
est infรฉrieure ร $`n`$. $`\mathrm{}`$
Un รฉlรฉment du revรชtement universel $`\stackrel{~}{\mathrm{Sp}(\mathrm{E},\omega )}`$ du groupe symplectique $`\mathrm{Sp}(\mathrm{E},\omega )`$ est la donnรฉe dโun chemin $`\gamma :[\mathrm{0,1}]\mathrm{Sp}(\mathrm{E},\omega )`$, tel que $`\gamma (0)=\mathrm{Id}`$, dรฉfini ร une homotopie fixant les extrรฉmitรฉs prรจs. On le note $`[\gamma ]\stackrel{~}{\mathrm{Sp}(\mathrm{E},\omega )}`$. On note $`\stackrel{~}{\mathrm{Id}}`$ lโรฉlรฉment dรฉfini par le chemin constant รฉgal ร $`\mathrm{Id}`$. Si $`L_0\mathrm{\Lambda }(E)`$, on note $`\phi _{L_0}([\gamma ])=\mathrm{\Delta }(det^2(\gamma _tL_0))`$. Si $`L_0`$ et $`L_1`$ sont deux lagrangiens on a
$$|\phi _{L_0}([\gamma ])\phi _{L_1}([\gamma ])|2n$$
(il suffit de considรฉrer lโapplication $`(s,t)[\mathrm{0,1}]^2det_{}^2(\gamma _tL_s)`$, oรน $`L_s`$ est un chemin de $`L_0`$ ร $`L_1`$ qui reste transverse ร un lagrangien donnรฉ, et dโappliquer la proposition prรฉcรฉdente ). De plus nous avons lโรฉgalitรฉ
$$\phi _{L_0}([\gamma ][\eta ])\phi _{L_0}([\gamma ])\phi _{L_0}([\eta ])=\phi _{\gamma _1L_0}([\eta ])\phi _{L_0}([\eta ]),$$
cette derniรจre quantitรฉ est bornรฉe par $`2n`$. Lโapplication $`\phi _{L_0}`$ est donc un quasi-morphisme sur le groupe $`\stackrel{~}{\mathrm{Sp}(\mathrm{E},\omega )}`$, dont lโhomogรฉnรฉisรฉ $`\mathrm{\Phi }`$ ne dรฉpend pas de $`L_0`$. Pour tout lagrangien $`L`$ on a :
$$\mathrm{\Phi }([\gamma ])=\mathrm{lim}_p\mathrm{}\frac{1}{p}\mathrm{\Delta }(det^2(\gamma _t^pL)).$$
###### Proposition 3.2.
Le quasi-morphisme $`\mathrm{\Phi }`$ est continu.
Preuve : on montre par rรฉcurrence sur $`k`$ lโinรฉgalitรฉ :
$$|\phi _{L_0}(x^{kp})k\phi _{L_0}(x^p)|2nk.$$
En divisant par $`kp`$ et en faisant tendre $`k`$ vers lโinfini, nous obtenons :
$$|\mathrm{\Phi }(x)\frac{1}{p}\phi _{L_0}(x^p)|\frac{2n}{p}.$$
La continuitรฉ de $`\phi _{L_0}`$ implique alors celle de $`\mathrm{\Phi }`$. $`\mathrm{}`$
Nous avons vu que le quasi-morphisme $`\mathrm{\Phi }:\stackrel{~}{\mathrm{Sp}(\mathrm{E},\omega )}`$ ne dรฉpend dโaucun choix de point base dans la lagrangienne. Il est รฉgalement indรฉpendant du choix de la structure presque-complexe. En effet, si $`T`$ est le gรฉnรฉrateur du groupe infini cyclique
$$\pi _1(\mathrm{Sp}(\mathrm{E},\omega ),\mathrm{Id})\stackrel{~}{\mathrm{Sp}(\mathrm{E},\omega )},$$
on vรฉrifie aisรฉment que $`\mathrm{\Phi }(T)=2`$. Ainsi les deux quasi-morphismes homogรจnes construits ร partir de deux structures presque-complexes $`J`$ et $`J^{}`$ distinctes, prennent la mรชme valeur sur lโรฉlรฉment $`T`$. Dโaprรจs un argument de Barge et Ghys , cela entraรฎne quโils sont รฉgaux. Indiquons finalement que lโon peut donner dโautres descriptions du quasi-morphisme $`\mathrm{\Phi }`$, qui permettent de le calculer effectivement .
Passons alors ร la construction du quasi-morphisme homogรจne $`\tau _{B,\nu }:\mathrm{\Gamma }_{B,\nu }`$. Nous choisissons une trivialisation symplectique du fibrรฉ tangent ร $`B`$. Si $`f\mathrm{\Gamma }_{B,\nu }`$, la diffรฉrentielle de $`f`$ โlueโ dans cette trivialisation est une application
$$\begin{array}{ccc}\hfill B& & \mathrm{Sp}(2n,)\hfill \\ \hfill x& & df(x).\hfill \end{array}$$
Un changement de trivialisation est donnรฉ par une application $`\theta :B\mathrm{Sp}(2n,)`$. La diffรฉrentielle de $`f`$ est alors changรฉe en $`x\theta (f(x))^1df(x)\theta (x)`$. Notons $`\stackrel{~}{\theta }:B\stackrel{~}{\mathrm{Sp}(2n,)}`$ un relevรฉ quelconque de $`\theta `$, et $`\stackrel{~}{df}:B\stackrel{~}{\mathrm{Sp}(2n,)}`$ lโunique relevรฉ de $`df`$ qui vaut $`\stackrel{~}{\mathrm{Id}}`$ hors dโun compact. Si $`f`$ et $`g`$ sont dans le groupe $`\mathrm{\Gamma }_{B,\nu }`$, nous avons :
$$\stackrel{~}{d(fg)}(x)=\stackrel{~}{df}(g(x))\stackrel{~}{dg}(x),$$
et donc
$$|\mathrm{\Phi }(\stackrel{~}{d(fg)}(x))\mathrm{\Phi }(\stackrel{~}{dg}(x))\mathrm{\Phi }(\stackrel{~}{df}(g(x)))|2n.$$
Lโapplication $`f_B\mathrm{\Phi }(\stackrel{~}{df})\nu ^n`$ est donc un quasi-morphisme sur le groupe $`\mathrm{\Gamma }_{B,\nu }`$. Nous allons vรฉrifier que son homogรฉnรฉisรฉ ne dรฉpend pas de la trivialisation symplectique choisie. Si lโon change de trivialisation, lโapplication $`\stackrel{~}{df}`$ est changรฉe en $`\stackrel{~}{\theta }^1f\stackrel{~}{df}\stackrel{~}{\theta }`$. Nous avons donc, puisque $`\mathrm{\Phi }`$ est homogรจne :
$$|\mathrm{\Phi }(\stackrel{~}{df^p}(x))\mathrm{\Phi }(\stackrel{~}{\theta }^1(f^p(x))\stackrel{~}{df^p}(x)\stackrel{~}{\theta }(x))|2n+|\mathrm{\Phi }(\stackrel{~}{\theta }^1(f^p(x))\stackrel{~}{\theta }(x))|.$$
Si le support de $`f`$ est contenu dans le compact $`K`$ de $`B`$, nous avons pour $`x`$ dans $`K`$, $`|\mathrm{\Phi }(\stackrel{~}{\theta }^1(f^p(x))\stackrel{~}{\theta }(x))|2n+2sup_K|\mathrm{\Phi }\stackrel{~}{\theta }|`$ ; si $`x`$ nโest pas dans $`K`$, la quantitรฉ $`\mathrm{\Phi }(\stackrel{~}{\theta }^1(f^p(x))\stackrel{~}{\theta })(x)`$ est nulle. Ainsi :
$$|\mathrm{\Phi }(\stackrel{~}{df^p}(x))\mathrm{\Phi }(\stackrel{~}{\theta }^1(f^p(x))\stackrel{~}{df^p}(x)\stackrel{~}{\theta }(x))|4n+2\mathrm{s}\mathrm{u}\mathrm{p}_K|\mathrm{\Phi }\stackrel{~}{\theta }|,$$
pour tout $`x`$ de $`B`$ et tout entier $`p`$. En divisant par $`p`$, en intรฉgrant, et en passant ร la limite, nous obtenons bien le rรฉsultat voulu : la quantitรฉ
$$\tau _{B,\nu }(f)=\mathrm{lim}_p\mathrm{}\frac{1}{p}_B\mathrm{\Phi }(\stackrel{~}{df^p})\nu ^n$$
ne dรฉpend pas de la trivialisation choisie. Lโapplication $`\tau _{B,\nu }:\mathrm{\Gamma }_{B,\nu }`$ est le quasi-morphisme homogรจne annoncรฉ.
### 3.2. Le quasi-morphisme $`๐`$
Notons dโabord que, puisque la classe $`[\omega ]`$ nโest jamais nulle, notre hypothรจse force $`c_1(V)0`$. Cela exclut par exemple les variรฉtรฉs symplectiques telles que les surfaces $`K3`$ ou les tores (on pourra consulter pour la construction dโun quasi-morphisme homogรจne sur le groupe $`\stackrel{~}{\mathrm{Ham}}(V,\omega )`$, pour les variรฉtรฉs de premiรจre classe de Chern nulle). Dans le cas oรน $`V`$ est ou bien $`^2`$ ou bien $`^1\times ^1`$ muni du produit de la forme symplectique standard par elle-mรชme, des rรฉsultats de M. Gromov assurent que le groupe fondamental de $`\mathrm{Ham}(V,\omega )`$ est fini. Tout quasi-morphisme homogรจne sur $`\stackrel{~}{\mathrm{Ham}}(V,\omega )`$ descend donc sur $`\mathrm{Ham}(V,\omega )`$. Notre construction fournit donc un nouvel exemple de quasi-morphisme homogรจne sur les groupes $`\mathrm{Ham}(^2,\omega _0)`$ et $`\mathrm{Ham}(^1\times ^1,\omega _0\times \omega _0)`$.
Fixons un $`S^1`$-fibrรฉ principal $`\pi :MV`$ de classe dโEuler รฉgale ร $`2c_1(V)`$. Notant $`X`$ le champ de vecteurs sur $`M`$ engendrรฉ par lโaction du cercle, on peut trouver une $`1`$-forme $`\alpha `$ sur $`M`$ telle que $`\alpha (X)=1`$ et $`d\alpha =\pi ^{}(s\omega )`$. Nous sommes dans la situation du paragraphe $`2.1`$. Fixons รฉgalement une structure presque-complexe $`J`$ sur $`V`$, compatible avec $`\omega `$. Le fibrรฉ vectoriel $`TV`$ devient alors un fibrรฉ hermitien, dont on peut choisir une trivialisation au-dessus dโun recouvrement $`\{U_\beta \}`$, avec des applications de transition $`g_{\beta \gamma }:U_\beta U_\gamma U(n)`$, ร valeurs dans le groupe des matrices unitaires de taille $`n`$. La famille dโapplications $`(det^2(g_{\beta \gamma }))`$, dรฉtermine un fibrรฉ en cercles $`E`$ au-dessus de $`V`$, qui est isomorphe ร $`M`$. Si lโon note $`\mathrm{\Lambda }(V)`$ le fibrรฉ en grassmannienne lagrangienne au-dessus de $`V`$, on a une application $`det^2:\mathrm{\Lambda }(V)E`$ qui nโest autre quโune version fibrรฉe de lโapplication dรฉjร rencontrรฉe dans le cas linรฉaire. Dans une trivialisation $`U_\gamma \times ^n`$, un รฉlรฉment $`L`$ de $`\mathrm{\Lambda }(V)`$ sโรฉcrit $`(x,u_\gamma (^n))`$, pour une matrice unitaire $`u_\gamma `$. On lui associe lโรฉlรฉment $`(x,det^2(u_\gamma ))`$ dans la trivialisation correspondante de $`E`$. En choisissant un isomorphisme entre $`E`$ et $`M`$ on obtient une application $`\phi :\mathrm{\Lambda }(V)M`$. Elle nโest bien sรปr pas unique, le choix de $`J`$ et lโisomorphisme entre $`E`$ et $`M`$ interviennent. Cependant, elle a la vertu suivante : en restriction ร chaque fibre, elle induit un isomorphisme entre les groupes fondamentaux. Une autre application $`\phi ^{}:\mathrm{\Lambda }(V)E`$ construite par le mรชme procรฉdรฉ serait donc de la forme
$$\phi ^{}(L)=\chi (\pi (L))e^{2i\pi \kappa (L)}\phi (L),$$
pour des application $`\chi :VS^1`$ et $`\kappa :\mathrm{\Lambda }(V)`$.
Passons ร la construction de notre dernier quasi-morphisme. On considรจre une isotopie hamiltonienne $`(f_t)`$ engendrรฉe par le champ de vecteurs $`Z_t`$ (avec $`\iota _{Z_t}\omega =dH_t`$, $`_VH_t\omega ^n=0`$). On note toujours $`\mathrm{\Theta }(f_t)`$ lโisotopie de $`M`$ engendrรฉe par le champ de vecteurs $`\widehat{Z_t}(H_t\pi )X`$. Si $`L\mathrm{\Lambda }(V)`$, les deux courbes
$$\phi (df_tL)\mathrm{et}\mathrm{\Theta }(f_t)(\phi (L))$$
dans $`M`$, sont issues du mรชme point et relรจvent la mรชme courbe de $`V`$. Bien que le fibrรฉ en cercles $`M`$ ne soit pas trivial, on peut se servir de la courbe $`\mathrm{\Theta }(f_t)(\phi (L))`$ comme dโune horizontale โle long du chemin $`f_t(\pi (L))`$โ pour mesurer le nombre de rotation de la courbe $`\phi (df_tL)`$. On peut donc รฉcrire $`\phi (df_tL)=e^{2i\pi \vartheta (t))}\mathrm{\Theta }(f_t)(\phi (L))`$ et dรฉfinir une fonction continue sur $`\mathrm{\Lambda }(V)`$ par $`\mathrm{angle}(L,\{f_t\})=\vartheta (1)\vartheta (0)`$. Elle satisfait la relation :
$$\mathrm{angle}(L,\{f_tg_tf_1\}=\mathrm{angle}(L,\{f_t\})+\mathrm{angle}(df_1L,\{g_t\}).$$
###### Proposition 3.3.
Pour toute paire de lagrangiens $`(L_0,L_1)`$ contenus dans la mรชme fibre de $`\mathrm{\Lambda }(V)V`$, et toute isotopie hamiltonienne $`\{f_t\}`$, nous avons :
$$|\mathrm{angle}(L_0,\{f_t\})\mathrm{angle}(L_1,\{f_t\})|2n.$$
Preuve : cโest une version fibrรฉe des rรฉsultats du paragraphe $`3.1`$, qui permettent de construire le quasi-morphisme homogรจne sur le revรชtement universel du groupe symplectique. $`\mathrm{}`$
Nous dรฉfinissons alors une fonction mesurable bornรฉe sur $`V`$ par : $`\mathrm{angle}(x,\{f_t\})=\mathrm{inf}_{L\mathrm{\Lambda }(V)_x}\mathrm{angle}(L,\{f_t\})`$. Elle satisfait :
$$|\mathrm{angle}(x,\{f_tg_tf_1\})\mathrm{angle}(x,\{f_t\})\mathrm{angle}(f_1(x),\{g_t\})|6n.$$
Lโapplication
$$\begin{array}{ccc}\hfill \stackrel{~}{\mathrm{Ham}}(V,\omega )& & \hfill \\ \hfill \{f_t\}& & _V\mathrm{angle}(,\{f_t\})\omega ^n\hfill \end{array}$$
est donc un quasi-morphisme. Si lโapplication $`\phi `$ est modifiรฉe en une application $`\phi ^{}`$, comme expliquรฉ ci-dessus, la fonction $`\mathrm{angle}`$ se trouve changรฉe en la fonction $`\mathrm{angle}^{}(x,\{f_t\})`$ รฉgale ร :
$$\mathrm{angle}(L,\{f_t\})+\kappa (df_1L)\kappa (L)+_0^1\beta (X_t)(f_t(x))๐t,$$
si $`\beta `$ dรฉsigne la $`1`$-forme fermรฉe $`d(\frac{ln\chi }{2i\pi })`$. Nous avons donc :
$$|\mathrm{angle}^{}(x,\{f_t\})\mathrm{angle}(x,\{f_t\})_0^1\beta (X_t)(f_t(x))๐t|4n+2\mathrm{s}\mathrm{u}\mathrm{p}_{\mathrm{\Lambda }(V)}|\kappa |.$$
Lโintรฉgrale $`_V_0^1\beta (X_t)๐t\omega ^n`$ รฉtant nulle, on a :
$$|_V\mathrm{angle}^{}(,\{f_t\})\omega ^n_V\mathrm{angle}(,\{f_t\})\omega ^n|(4n+2\mathrm{s}\mathrm{u}\mathrm{p}_{\mathrm{\Lambda }(V)}|\kappa |)\mathrm{vol}(V).$$
Lโhomogรฉnรฉisรฉ
$$๐(\{f_t\})=\mathrm{lim}_p\mathrm{}\frac{1}{p}_V\mathrm{angle}(,\{f_t\}^p)\omega ^n$$
de notre quasi-morphisme ne dรฉpend donc pas du choix de lโapplication $`\phi `$, et donc pas de $`J`$. Il ne dรฉpend pas non plus du choix de la $`1`$-forme $`\alpha `$ (vรฉrifiant $`\alpha (X)=1`$ et $`d\alpha =\pi ^{}(s\omega )`$). Nous calculons maintenant sa restriction sur les isotopies ร support dans une boule $`i:BV`$.
###### Proposition 3.4.
Si lโisotopie $`\{f_t\}`$ est ร support dans $`B`$, nous avons
$$๐(\{f_t\})=\tau _{B,\omega }(f_1)+s๐๐ฉ_B(f_1).$$
Preuve : on fixe une trivialisation unitaire du fibrรฉ tangent au-dessus de $`B`$, et une trivialisation du fibrรฉ en cercles $`M`$ au-dessus de $`B`$. On note $`\lambda `$ la primitive de $`\omega `$ sur $`B`$ telle que $`\alpha =d(\frac{lnz}{2i\pi })+s\lambda `$ dans cette trivialisation ($`z`$ dรฉsigne la coordonnรฉe sur le cercle). Lโapplication $`\phi `$ lue dans cette trivialisation est de la forme :
$$\begin{array}{ccc}\hfill B\times \mathrm{\Lambda }(^{2n})& & B\times S^1\hfill \\ \hfill (x,L)& & (x,e^{2i\pi \kappa (x,L)}det_^n^2(L))\hfill \end{array}$$
pour une application $`\kappa :B\times \mathrm{\Lambda }(^n)`$. On considรจre alors une isotopie hamiltonienne $`(f_t)`$ ร support contenu dans un compact $`K`$ de $`B`$, engendrรฉe par le champ de vecteurs $`Z_t`$ (avec $`\iota _{Z_t}\omega =d\stackrel{~}{H}_t`$, oรน $`\stackrel{~}{H}_t`$ est dโintรฉgrale nulle sur $`V`$, constante hors de $`B`$). Nous avons dโune part :
$$\phi (f_t(x),df_tL)=(f_t(x),e^{2i\pi \kappa (f_t(x),df_tL)}det_^n^2(df_tL))$$
et dโautre part :
$$\mathrm{\Theta }(f_t)(\phi (x,L))=(f_t(x),e^{2i\pi _0^t(\lambda (Z_u)+s\stackrel{~}{H}_u)(f_u(x))๐u}e^{2i\pi \kappa (x,L)}det_^n^2(L)).$$
La valeur de $`\mathrm{angle}(L,\{f_t\})`$ est donc :
$$\mathrm{\Delta }(det_^n^2(df_t(x)L))+_0^1(s\lambda (X_t)+sH_t)(f_t(x))๐t+\kappa (f_1(x),df_1(x)L)\kappa (x,L).$$
Notons $`C`$ le maximum de $`|\kappa |`$ sur $`\mathrm{supp}(f)\times \mathrm{\Lambda }(^{2n})`$. Nous obtenons
$$|\mathrm{angle}(x,\{f_t\})\mathrm{\Delta }(det_^n^2(df_t(x)L))s_0^1(\lambda (X_t)+\stackrel{~}{H}_t)(f_t(x))๐t|2C+2n.$$
Hors de la boule $`B`$ la fonction $`\mathrm{angle}(x,\{f_t\})`$ est รฉgale ร $`s_0^1\stackrel{~}{H}_t(f_t(x))๐t`$. En tenant compte de lโinรฉgalitรฉ $`|\mathrm{\Delta }(det_^n^2(df_t(x)L))\mathrm{\Phi }(\{df_t(x)\})|2n`$, nous obtenons :
$$|_V\mathrm{angle}(,\{f_t\}\nu ^ns_0^1_B\lambda (X_t)dt\nu ^n_B\mathrm{\Phi }(\{df_t(x)\})\nu ^n|(2C+4n)\mathrm{vol}(B).$$
La mรชme estimation reste vraie pour les itรฉrรฉs de $`f`$ car leur support est contenu dans celui de $`f`$. Nous obtenons donc : $`๐(\{f_t\})=s_0^1_B\lambda (X_t)๐t\omega ^n+\tau _{B,\nu }(f_1)=s๐๐ฉ(f_1)+\tau _{B,\nu }(f_1)`$. $`\mathrm{}`$
Remerciements. รtienne Ghys mโa proposรฉ de rรฉflรฉchir ร ce sujet et mโa encouragรฉ; je le remercie pour tout cela. Je tiens รฉgalement ร remercier Emmanuel Giroux et Bruno Sรฉvennec pour de nombreuses discussions. Enfin, je tiens ร remercier vivement Leonid Polterovich qui mโa signalรฉ une erreur dans une version antรฉrieure du thรฉorรจme $`2`$.
Pierre Py
Unitรฉ de Mathรฉmatiques Pures et Appliquรฉes
UMR 5669 CNRS
รcole Normale Supรฉrieure de Lyon
46, Allรฉe dโItalie
69364 Lyon Cedex 07
FRANCE
Pierre.Py@umpa.ens-lyon.fr
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# Quantum incompressibility and Razumov Stroganov type conjectures
## 1 Introduction
This paper is aimed at establishing a correspondence between the deformation of certain wave functions of the Hall effect and polynomial representations of the Temperley and Lieb (T.L.) algebra.
This work originates from an attempt to understand the conjecture of A.V. Razumov and Y.G. Stroganov , and some partial results towards its proof by P.Di Francesco and P. Zinn-Justin .
We consider the analogue of spin singlet wave functions of the Hall effect when one deforms the permutations into the braid group. This amounts to analyze some simple representations of the T.L. algebra on a space of polynomials in $`N_e`$ variables where $`N_e`$ is the number of electrons. The relation with the Hall effect arises when we require certain incompressibility properties.
One of the wave functions we consider here is the Halperin wave function for a system of spin one half electrons at filling factor two. When the deformation parameter $`q`$ is a third root of unity, the braid group representation degenerates into a trivial representation. In this way, we obtain a proof alternative to, and apparently simpler than, that given in of the equality between the sum of the components of the transfer matrix eigenvector and the six vertex model partition function with domain wall boundary conditions .
Another wave function we consider is the Haldane Rezayi wave function <sup>1</sup><sup>1</sup>1More precisely a minor modification of it considered in . describing a system of electrons of spin one half at filling factor one. This wave function is a permanent, and its deformation is described in terms of Gaudinโs determinants . When $`q`$ is a third root of unity, it degenerates to the square of the six vertex model partition function.
In a separate publication , we shall consider the Moore Read wave function describing spinless bosons at filling factor one . Its deformation involves an extension of the braid group known as the Birman-Wenzl algebra which can be represented on a polynomial space similarly to the cases presented here. In some appropriate limit, the representation degenerates and the wave function coincides with the transfer matrix eigenvector considered in related to the conjecture of J.De Gier and B. Nienhuis .
In general, when a Quantum Hall Effect wave function is discovered, it is soon after observed experimentally. We argue here, that as a bonus, Quantum Hall Effect wave functions and their deformations yield nice mathematical objects. Moreover, all these objects seem to be in relation with striking conjectures emanating from the six vertex model.
Since the permutation group relevant in the quantum Hall effect is technically simpler than the braid group case, let us for pedagogical reasons explain why finding a wave function turns out to be a useful tool to obtain a polynomial representations of the permutation algebra. Essentially, the rest of the paper extends the idea presented here to the braid group case.
We consider electrons in a strong magnetic field projected in the lowest Landau level. In a specific gauge the orbital wave functions are given by:
$`\psi _n(z)={\displaystyle \frac{z^n}{\sqrt{n!}}}e^{\frac{z\overline{z}}{4l^2}},`$ (1)
where $`z=x+iy`$ is the coordinate of the electron, and $`l`$ the magnetic length defines the length scale related to the strength of the magnetic field. These orbitals are shells of radius $`\sqrt{2n}l`$ occupying an area $`2\pi l^2`$. Each orbital $`n`$ is represented by a monomial $`z^n`$.
The quantum Hall effect ground state $`\mathrm{\Psi }`$ is obtained by combining these individual orbitals into a manybody wave function. A monomial $`z_1^{\lambda _1}\mathrm{}z_{N_e}^{\lambda _{N_e}}`$ describes a configuration where the electron $`j`$ occupies the orbital $`\lambda _j`$. The wave function is a linear combinations of such monomials. The effect of the interactions is to impose some vanishing properties when electrons are in contact: $`\mathrm{\Psi }(z_iz_j)^m`$ with $`m`$ an integer when $`z_iz_j0`$.
The physical properties are mainly characterized by the filling factor $`\nu `$ which is the number of electrons per unit cell of area $`2\pi l^2`$. When the filling factor is equal to $`\nu `$, the accessible orbitals and thus the maximal degree in each variable is bounded by $`\nu ^1N_e`$. On the other hand, the effect of the interactions ($`m`$) is to force the electrons to occupy more space, thus to occupy higher orbitals and and has the effect of increasing the degree. The problem is thus to obtain wave functions with the maximal possible filling factor (equivalently the lowest degree in each variable) compatible with the vanishing properties imposed by the interactions.
Once such a wave function is obtained, it is the nondegenerate lowest energy state of a Hamiltonian invariant under the permutations, thus we know that it is left invariant under the permutations. By disentangling the coordinate part from the spin part, we obtain an irreducible representation of the permutation algebra acting on polynomials.
Let us illustrate this point in the case of the Halperin wave function which describes a system of spin one half electrons at filling factor two. There are no interactions between the electrons, but due to the Pauli principle, the wave function must vanish when two electrons of the same spin come into contact. Each independent orbital can be occupied with two electrons of opposite spin, which is why the maximal filling factor is equal to two.
An equivalent way to impose the constraint is to require that any linear combination of the spin components of the wave function vanishes when three electrons come into contact. The reason for this is that two of the electrons involved will necessary have the same spin. When this constraint is taken into account with the minimal degree hypotheses, one obtains a space of polynomials which can be recombined with the spin components into a wave function changing sign under the permutations. Thus we know a priori that the spatial part of the wave function carries an irreducible representation of the permutation algebra dual to that of the spins. This is precisely by generalizing this argument to the braid group case that we obtain the representations of the T.L. algebra mentioned above.
In the permutation group case case, the components have the simple structure of a product of two Slater determinants grouping together the electrons with the same spin and one does not require to recourse to this machinery.
Let us now briefly indicate why the Halperin wave function may have something to do with the eigenvector of a transfer matrix in the link pattern formulation . The wave function is a spin singlet, and the spin components can best be described in a resonating valance bond (RVB) picture as follows: The labels of the electrons are disposed cyclically around a circle and are connected by a link when two electrons form a spin singlet. Links are not allowed to cross in order to avoid overcounting states. These RVB states coincide with the link patterns of . Thus, the Halperin wave function as the eigenvector of the transfer matrix develops on a basis of link patterns. By deforming the permutation action on link patterns into a T.L algebra action, one is forced to deform accordingly the polynomial representation so as to insure the invariance of the total wave function. When $`q`$ is a third root of unity, this property is shared by the transfer matrix eigenvector and allows to identify the two.
In the braid group case, the situation is technically more involved than for the permutations. Nevertheless, the minimal degree hypothesis combined with some annulation constraint satisfied by linear combination of the spin components yields a wave function with the correct invariance properties. A major difference with the Hall effect is that the cancelation no longer occurs at coincident points, but at points shifted proportionally to the deformation parameter $`q`$. Typically, we require that for three arbitrary electron labels $`i<j<k`$ ordered cyclically, the wave function vanishes when the corresponding coordinates take the values $`z,q^2z,q^4z`$.
One is also led to study the affine extension in order to impose cyclic invariance properties which are tautologically satisfied with the permutations. While defined in a natural way on the link patterns, the cyclic properties require to introduce a shift parameter $`s`$ when we identify the coordinate $`i+N_e`$ with the coordinate $`i`$: $`z_{i+N_e}=sz_i`$. When this shift parameter is related in a specific way to the braid group deformation parameter, the generalized statistics properties can be established coherently. Here, $`z_{i+N_e}=q^6z_i,`$ but the same annulation property can also be satisfied with $`s`$ not related to $`q`$, and this can be achieved at the price of doubling the degree and enlarging the algebra .
In the the Haldane Rezayi case, , the interactions are such that the wave function must vanish as the square of the distance when electrons of the same spin come into contact. For the same reason as before, this amounts to impose that any linear combination of its spin components vanishes as the square of the distance when three electrons come into contact. This wave function is a permanent, and its deformation is described in terms of Gaudinโs determinants . It degenerates to the square of the six vertex model partition function when the deformation parameter is a third root of unity.
The paper is organized as follows. In section 2, we recall some properties about Hecke algebras and their polynomial representations. Section 3 introduces the T.L. algebra representation used here. Section 4 is the core of the paper where we work out the deformed Hall effect wave functions.
We have attempted to be self contained, but in order not to overload the text with technicalities, we have relegated most of the proofs to appendices to which we refer when it is useful.
## 2 Hecke Algebra.
In this section, we recall some known facts about the Hecke and Temperley and Lieb algebras .
The Braid group algebra is generated by the braid group generators $`t_1,t_2,\mathrm{},t_{n1}`$, obeying the braid relations:
$`t_it_{i+1}t_i`$ $`=`$ $`t_{i+1}t_it_{i+1}`$ (2)
$`t_it_j`$ $`=`$ $`t_jt_i\mathrm{if}|ij|>1,`$ (3)
for $`1in1`$. It can be convenient to use the notation $`t_{ii+1}`$ instead of $`t_i`$, and we will use it when necessary. The Hecke algebra is the quotient of the Braid group algebra by the relations:
$`(t_iq)(t_i+{\displaystyle \frac{1}{q}})=0,`$ (4)
It can also be defined using the projectors $`e_i=t_iq`$ obeying the relations:
$`e_i^2`$ $`=`$ $`\tau e_i,`$ (5)
$`e_ie_j`$ $`=`$ $`e_je_i\mathrm{if}|ij|>1`$ (6)
$`e_ie_{i+1}e_ie_i`$ $`=`$ $`e_{i+1}e_ie_{i+1}e_{i+1}.`$ (7)
where we set $`\tau =(q+q^1)`$. The Temperley-Lieb (T.L.) algebra $`๐_n`$ is the quotient of Hecke algebra by the relations:
$`e_ie_{i+1}e_ie_i=e_{i+1}e_ie_{i+1}e_{i+1}=0.`$ (8)
In $`๐_n`$, a trace can be defined as:
$`\mathrm{tr}(xe_p)=\tau ^1\mathrm{tr}(x)x๐_p.`$ (9)
The affine Hecke algebra, , is an extension of the Hecke algebra (4) by generators $`y_i,1in`$ obeying the following relations:
$`a)`$ $`y_iy_j=y_jy_i`$ (10)
$`b)`$ $`t_iy_j=y_jt_i\mathrm{if}ji,i+1`$ (11)
$`c)`$ $`t_iy_{i+1}=y_it_i^1\mathrm{if}in1.`$ (12)
In (D.2), we indicate why (12c) is natural from the Yang-Baxter algebra point of view.
This algebra can be endowed with two possible involutions: $`e_i^{}=e_i,y_i^{}=y_i^{\pm 1}`$, $`q^{}=q^{\pm 1}`$.
The symmetric polynomials in the $`y_i`$ are central elements.
We define the affine T.L. algebra $`๐_n^{}`$ as the extension of the T.L. algebra (8) by the generators $`y_i`$.
### 2.1 Yangโs realization of the Affine relations.
The commutation relations of the affine generators $`y_i`$ become simpler to understand if we assume that we have a representation of the permutations $`k_{ij}`$ acting in the natural way on the indices. Let us introduce the operators $`x_{ij}=t_{ij}k_{ij}`$ for $`i<j`$ and $`x_{ji}=x_{ij}^1`$. These operators obey the Yangโs relations:
$`x_{ij}x_{ji}`$ $`=`$ $`1,`$ (13)
$`x_{ij}x_{kl}`$ $`=`$ $`x_{kl}x_{ij}ijkl,`$ (14)
$`x_{ij}x_{ik}x_{jk}`$ $`=`$ $`x_{ik}x_{jk}x_{ij}.`$ (15)
We also assume that we have commuting operators $`s_i`$ such that $`s_is_jx_{ij}=x_{ij}s_is_j`$.
Using (15), one verifies that the operators defined as :
$`y_1`$ $`=`$ $`x_{12}x_{13}\mathrm{}x_{1n}s_1`$ (16)
$`y_2`$ $`=`$ $`x_{23}x_{24}\mathrm{}x_{2n}s_2x_{21}`$ (17)
$`y_n`$ $`=`$ $`s_nx_{n1}x_{n2}\mathrm{}x_{nn1}.`$ (18)
commute. Indeed, they coincide with the scattering matrices of Yang . (12b) follows directly from (15) once we substitute $`t_i=x_{ii+1}k_{ii+1}`$. (12c) is a direct consequence of the definition (18) of $`y_i`$.
Gathering the permutation operators $`k_{ij}`$ together, we can obtain another presentation of the $`y_i`$. Let us introduce the cyclic operator:
$`\sigma =k_{n1n}\mathrm{}k_{23}k_{12}s_1.`$ (19)
Then we have:
$`y_1`$ $`=`$ $`t_1t_2\mathrm{}t_{n1}\sigma `$ (20)
$`y_2`$ $`=`$ $`t_1^1y_1t_1^1`$ (21)
$`y_n`$ $`=`$ $`t_{n1}^1y_{n1}t_{n1}^1.`$ (22)
We can define an additional generator to the $`t_i`$: $`t_n=\sigma t_1\sigma ^1`$, which makes the relations (3) become cyclic. One has:
$`\sigma t_i=t_{i1}\sigma .`$ (23)
So that the affine Hecke (or T.L.) algebra is generated by the generators $`t_i`$ and the cyclic operator $`\sigma `$ obeying (23) and does not require a representation of the permutations. $`\sigma ^n`$ is a central element which can be set equal to one, and $`\sigma ^{}=\sigma ^1`$ if we take $`t^{}=t^1`$.
Given the Hecke algebra, there is a simple realization of the affine Hecke algebra which consists in taking $`y_1=1`$. Then, $`\sigma `$ is defined as:
$`\sigma =t_{n1}^1\mathrm{}t_1^1.`$ (24)
Using the braid relations, one sees that $`\sigma t_i=t_{i1}\sigma `$ for $`i>1`$, and one can define $`t_n`$ by $`t_n=\sigma t_1\sigma ^1`$. Using the braid relations again, one gets $`\sigma t_n=t_{n1}\sigma `$. This defines an operator $`\sigma `$ which allows to construct the affine generators with (22).
### 2.2 Polynomial representations.
We consider a space of polynomials in $`n`$ variables $`z_i,`$ constructed as a linear combinations of monomials: $`z^\mu =z_1^{\mu _1}z_2^{\mu _2}\mathrm{}z_n^{\mu _n}`$ with a total degree $`|\mu |=\mu _i`$ fixed. There is a natural action of the permutations and of the operators $`s_i`$ on this space defined by:
$`\overline{\psi }(z_1,..z_i..z_j..,z_n)k_{ij}`$ $`=`$ $`\overline{\psi }(z_1,..z_j..z_i..,z_n),`$ (25)
$`\overline{\psi }(z_1,..,z_i..,z_n)s_i`$ $`=`$ $`c\overline{\psi }(z_1..,sz_i..,z_n).`$ (26)
It is convenient to consider the polynomials in an infinite set of variables $`z_i,i๐ต,`$ with the identification: $`z_{i+n}=sz_i`$. The operator $`\overline{\sigma }`$ (19) takes the form:
$`\overline{\psi }\overline{\sigma }(z_i)=c\overline{\psi }(z_{i+1}).`$ (27)
The condition $`\sigma ^n=1`$ imposes the relation $`c^ns^{|\mu |}=1`$.
As shown in (D.1), it is straightforward to derive the following representation of the Hecke relations (3,4):
$`\overline{t}_{ij}=q^1+(1k_{ij}){\displaystyle \frac{qz_iq^1z_j}{z_iz_j}}.`$ (28)
In this way we obtain a representation of the affine Hecke algebra acting on homogenous polynomials of a given total degree.
The operators $`x_{ij}`$ take the form:
$`x_{ij}=q^1+(qq^1)(1k_{ij}){\displaystyle \frac{z_j}{z_iz_j}}.`$ (29)
In the appendix D.3, we show that there is a natural order on the monomial basis, $`z^\mu `$, for which the operators $`x_{ij}`$, and hence the $`y_i`$ are realized as lower triangular matrices.
The operator $`y=y_1+\mathrm{}+y_n`$ can be seen to commute with the Hecke generators. It is therefore equal to a constant in an irreducible representation. Its eigenvalue evaluated on the highest weight polynomial $`P_\lambda `$ thus characterizes the representation. It is given by:
$`y_\lambda =c(q)^{1n}(s^{\lambda _1}+s^{\lambda _2}q^2+\mathrm{}+s^{\lambda _n}q^{2(n1)}).`$ (30)
If $`\lambda ^{}`$ is a permutation of the partition $`\lambda `$, we say that $`z^\lambda ^{}`$ is of degree $`\lambda `$. In this paper, we are mainly concerned with the monomials $`z^\lambda ^{}`$, of degree:
$`\lambda =({\displaystyle \frac{n}{2}}1,{\displaystyle \frac{n}{2}}1,{\displaystyle \frac{n}{2}}2,{\displaystyle \frac{n}{2}}2,\mathrm{},0,0),`$ (31)
and of total degree $`|\lambda |=\frac{n}{2}(\frac{n}{2}1)`$.
We will consider the subclass $`\lambda _\pi `$ of permutations of $`\lambda `$ which are smaller than $`\lambda `$ for the order introduced in the appendix D.3. According to the analysis made there, $`\lambda _\pi `$ can be obtained from $`\lambda `$ by a sequence of permutations $`(\lambda _i^{},\lambda _{i+1}^{})(\lambda _{i+1}^{},\lambda _i^{})`$ with $`\lambda _i^{}>\lambda _{i+1}^{}`$. The only monomials which can be obtained that way are indexed by the standard Young tableaus with two columns of $`\frac{n}{2}`$ boxes:
$`z^{\lambda _\pi }=(z_{\mu _1}z_{\nu _1})^{\frac{n}{2}}(z_{\mu _2}z_{\nu _2})^{\frac{n}{2}1}\mathrm{}(z_{\mu _{\frac{n}{2}}}z_{\nu _{\frac{n}{2}}})^0,`$ (32)
with $`\mu _1>\mu _2>\mathrm{}>\mu _{\frac{n}{2}}`$, $`\nu _1>\nu _2>\mathrm{}>\nu _{\frac{n}{2}}`$, and $`\mu _i>\nu _i`$. To simplify notations, we denote these monomials by $`z^\pi `$ instead of $`z^{\lambda _\pi }`$.
We identify the standard Young tableaus with the paths $`\pi =[h_i]`$ introduced in the appendix A: $`h_0=h_n=0`$, $`h_i0`$ and $`h_{i+1}h_i=\pm 1`$. These paths are obtained using the rule: $`h_ih_{i1}=1`$ if $`i\{\mu _j\}`$, and $`h_ih_{i1}=1`$ if $`i\{\nu _j\}`$. For the paths, we use the order $`\pi \pi ^{}`$, if $`[h_i][h_i^{}]i,`$ which coincides with the reverse order for the monomials: $`z^\pi z^\pi ^{}`$.
This identification is illustrated in figure 1.
## 3 Representation of the affine T.L. algebra on words.
For $`n`$ even, there is a simple representation $`(_n)`$ of the T.L. algebra $`๐_n`$ obtained as follows. One considers the left action of $`๐_n`$ on the space $`๐_n\alpha `$ where $`\alpha `$ is the minimal projector $`\alpha =e_1e_3\mathrm{}e_{n1}`$. A basis of this space is given by reduced monomial words in the $`e_i`$. The elements of this basis can be put into correspondence with paths or link patterns. In the appendix A we exhibit a basis of reduced words and we define an order relation on the reduced words.
A scalar product can be defined as:
$`\pi ^{}\pi ^{}=\pi |\pi ^{}\alpha ,`$ (33)
where $`e_i^{}=e_i`$ and the involution reverses the order of the letters. In the link-pattern representation, this scalar product is given by: $`\tau ^l`$ where $`l`$ is the number of loops one gets by stacking the link patterns of $`\pi `$ and $`\pi ^{}`$ on top of each other. If $`\tau =(q+q^1)`$ with $`q`$ not a root of one, this scalar product is positively definite . For this scalar product the T.L generators $`e_i`$ are by construction hermitian.
To obtain the affine algebra representation, let us define as in (24) the cyclic operator:
$`\sigma =q^{\frac{n}{2}2}t_{n1}^1\mathrm{}t_1^1,`$ (34)
where the normalization is such that in the link-pattern representation, $`\sigma `$ acts by cyclicly permuting the indices $`ii1`$ (see appendix A.4). One can define an additional generator, $`e_{n1}=\sigma e_{12}\sigma ^1`$, which acts in the same way as $`e_{ii+1}`$ with the two indices $`1,n`$. The affine generators are constructed using (22) with $`y_1=q^{\frac{n}{2}2}`$.
In the appendix A.3, we show that the operators $`y_i`$ are realized as triangular matrices in $`_n`$, they are hermitian for the choice $`q=q^{}`$. Their sum $`y=_iy_i`$ is constant with a value given by:
$`y=(q+q^1){\displaystyle \frac{q^{\frac{n}{2}}q^{\frac{n}{2}}}{qq^1}}.`$ (35)
There is an imbedding of $`_{n2}`$ into $`_n`$ given by $`\pi \pi e_1`$ and a projection $`E`$ from $`_n`$ to $`_{n2}`$ given by:
$`e_1\pi =\tau E(\pi )e_1.`$ (36)
This projection is both triangular and hermitian.
In A.2, we identify $`_n`$ with $`๐_{\frac{n}{2}}`$. This allows us to interpret the projection $`E`$ as a conditional expectation value of $`๐_{\frac{n}{2}}๐_{\frac{n}{2}1}`$ .
## 4 q-deformed Quantum Hall Effect wave functions.
### 4.1 Statement of the Problem.
Let us consider a vector $`\mathrm{\Psi }`$:
$`\mathrm{\Psi }={\displaystyle \underset{\pi }{}}\pi F_\pi (z_i),`$ (37)
constructed in the following way. The vectors $`\pi `$ are the basis vectors of $`_n`$ on which the T.L. algebra acts to the left. $`F_\pi `$ are homogeneous polynomials in the variables $`z_1,z_2,\mathrm{},z_n`$ ($`n`$ is even). The polynomial coefficients of $`\mathrm{\Psi }`$ carry a representation of the affine Hecke algebra generated by the operators $`\overline{t}_i`$ and $`\overline{\sigma }`$ acting to the right. The problem is to determine the coefficients $`F_\pi `$ in such a way that both actions give the same result on the vector $`\mathrm{\Psi }`$:
$`\mathrm{\Psi }\overline{t}_i`$ $`=`$ $`t_i\mathrm{\Psi }`$ (38)
$`\mathrm{\Psi }\overline{\sigma }`$ $`=`$ $`\sigma \mathrm{\Psi },`$ (39)
The first of these relations is equivalent to the more familiar relation (83) derived in D.1.
Said differently, we look for a dual action of the affine T.L. algebra acting on polynomials. Unless we specify it, we address this problem for a generic value of the parameter $`q`$, not a root of unity, for which the T.L. algebra is semisimple .
### 4.2 Module $`_n`$.
The dual representation of $`_n`$ is obtained by acting with the T.L. generators on the dual $`F_\omega `$ of the highest vector $`\omega _n`$. $`\omega `$ is given by the sequence $`(a_{2p+1}=p+1)`$ in the characterization of words we use in the appendix A and is fully characterized by the property that it can be written $`\omega =e_i\pi `$ only for $`i=\frac{n}{2}`$. The dual vector $`F_\omega `$ must therefore be annihilated by all the $`e_i`$ with $`in/2`$. We realize the module $`_n`$ upon acting on $`F_\omega `$ with the generators $`e_i`$ for $`1in1`$. We define:
$`_n=\mathrm{๐๐๐}\{\overline{\psi }=F_\omega \psi \},`$ (40)
where we denote with a bar $`\overline{\psi }`$ the result of the action of the monomial $`\psi `$ on $`F_\omega `$. Thus we have $`\overline{1}=F_\omega `$. In the appendix C.1, we show that $`_n`$ defined in this way is a module over the T.L. algebra as long as the $`e_i`$ obey the Hecke relations (7). In other words, the projectors $`U_{i,i+1}^{}=e_ie_{i+1}e_ie_i`$ are null in $`_n`$. This formal module is however not isomorphic to $`_n`$ unless $`F_\omega `$ obeys some supplementary condition (76). Here, we construct a representation of the T.L. algebra by identifying a state $`F_\omega `$ dual to $`\omega `$ and satisfying the condition (76).
The expression of the T.L. generators $`e_i=t_iq`$ for $`1in1`$ follows from (28):
$`e_i`$ $`=`$ $`{\displaystyle \frac{qz_{i+1}q^1z_i}{z_{i+1}z_i}}(1+k_{ii+1})`$ (41)
$`e_i\tau `$ $`=`$ $`(1k_{ii+1}){\displaystyle \frac{qz_iq^1z_{i+1}}{z_iz_{i+1}}}.`$ (42)
The effect of $`e_i`$ and $`\tau e_i`$ is to split a polynomial $`\overline{\psi }`$ into two polynomials belonging to $`_n`$, $`\overline{\psi }=S_1+(qz_iq^1z_{i+1})S_2`$, where both $`S_1`$ and $`S_2`$ are symmetrical under the exchange of $`z_i`$ and $`z_{i+1}`$. This decomposition is unique and characterizes the projector $`e_i`$.
It can be convenient to distinguish the representation on $`_n`$ from its dual on $`_n`$. When this is the case, we denote $`\overline{e}_i`$ the dual projectors which act on polynomials.
One verifies that:
$`\mathrm{\Delta }_m^q(z_1,\mathrm{},z_p)={\displaystyle \underset{1i<jm}{}}(qz_iq^1z_j)`$ (43)
is annihilated by all the $`e_i,1im1`$, and this defines $`\mathrm{\Delta }_m`$ up to a product by a symmetric polynomial in $`z_1,\mathrm{},z_m.`$ Therefore, the minimal degree polynomial candidate for $`F_\omega `$ is:
$`F_\omega =\mathrm{\Delta }_{\frac{n}{2}}^q(z_1,\mathrm{},z_{\frac{n}{2}})\mathrm{\Delta }_{\frac{n}{2}}^q(z_{\frac{n}{2}+1},\mathrm{},z_n).`$ (44)
This polynomial cannot be q-antisymmetrized over $`\frac{n}{2}+1`$ variables, since the result would have a degree at least $`\frac{n}{2}`$ in $`z_1`$, and this is the content of the condition (76). Thus, $`_n`$ is a simple module which can be identified with $`_n`$. This representation is characterized by its Young diagram $`(2^{\frac{n}{2}})`$ having two columns of length $`\frac{n}{2}`$. Its dimension is given by the Catalan number $`C_n=\left(\genfrac{}{}{0pt}{}{n}{\frac{n}{2}}\right)\left(\genfrac{}{}{0pt}{}{n}{\frac{n}{2}1}\right)`$.
If we denote $`\pi |_\omega `$ the coefficient of $`\omega `$ in the reduced expression of $`\pi `$, we can identify the polynomials $`\overline{\psi }_n`$ with the dual of $`_n`$ through the relation: $`\overline{\psi }(\pi )=\psi \pi |_\omega `$.
We can also introduce the dual basis $`F_\pi `$ defined by its action on reduced words:
$`F_\pi (\pi ^{})=\delta _{\pi ,\pi ^{}}.`$ (45)
Let $`\pi _\psi `$ be the complementary word of $`\psi `$ (defined in appendix C) such that one can write $`\psi \pi _\psi =\omega `$ without reducing the expression. One has $`\psi \pi _\psi |_\omega =1`$ and $`\psi \pi |_\omega =0`$ if $`\pi <\pi _\psi `$ . Expanding $`\overline{\psi }`$ on the basis $`F_\pi `$, we get: $`\overline{\psi }=_{\pi \pi _\psi }\psi \pi |_\omega F_\pi `$, and by inverting the triangular system, we can obtain the expression of $`F_\pi `$.
Let us verify that the highest monomial of $`\overline{\psi }`$, and thus of $`F_{\pi _\psi }`$ as well, is proportional to $`z^{\pi _\psi }`$. We show this by recursion. It is true for $`\psi =1`$: $`\overline{1}=F_\omega `$ , $`\pi _1=\omega `$ and $`z^\omega `$ is the highest monomial of $`F_\omega `$. We assume that the property is true for $`\psi ^{}<\psi `$. If $`\psi 1`$, one can write $`\overline{\psi }=\overline{\psi }^{}e_i`$ with $`\overline{\psi }^{}<\overline{\psi }`$, and we have $`\pi _\psi ^{}=e_i\pi _\psi `$ with $`\pi _\psi <\pi _\psi ^{}`$.
Then, according to the recursion hypothesis, the highest monomial of $`\overline{\psi }^{}`$ is $`z^{\pi _\psi ^{}}`$ which contains the factor $`z_i^mz_{i+1}^n`$ with $`n>m`$. Since $`z_i^mz_{i+1}^n\overline{e}_i=qz_i^nz_{i+1}^m+\mathrm{lower}\mathrm{monomials},`$ the highest monomial of $`\overline{\psi }=\overline{\psi }^{}\overline{e}_i`$ is $`z^{\pi _\psi }.`$
We also obtain the normalization coefficient of $`z^\pi `$ up to a global factor: $`F_\pi =c_\pi z^\pi +\mathrm{lower}\mathrm{monomials},`$ with $`c_\pi =(\frac{1}{q})^{l_\pi }`$, and $`l_\pi `$ is the number of letters $`e_i`$ entering the reduced expression of $`\pi `$.
### 4.3 Module $`_n^{}`$.
We now consider a larger module $`_n^{}_n`$ by letting the operator $`\overline{\sigma }`$ defined in (27) act on the polynomials. We will put some constraint on the parameter $`s`$ (which characterizes $`\overline{\sigma }`$) to have $`_n^{}=_n`$. We consider the simple case $`n=4`$ in the appendix B and we obtain $`s=q^6`$ which is the general case as we show here.
Let us assume that $`_n^{}=_n`$ and see what constraints $`s`$ must satisfy to identify $`\overline{\sigma }`$ defined by its action on polynomials (27) with $`\sigma `$ defined in terms of generators (24).
We observe that $`\sigma ^1\omega _n=e_1\omega _{n2}`$, where $`\omega _{n2}`$ is the highest state in $`_{n2}`$. This can easily be verified in the link pattern representation. Thus, we must have:
$`\sigma E(\pi )e_1|_{\omega _n}=E(\pi )e_1|_{\sigma ^1\omega _n}=E(\pi )|_{\omega _{n2}}.`$ (46)
Let us consider the dual to the projection $`E`$, $`E^{}`$ from $`_n_{n2}`$ defined as $`\overline{\psi }e_1=\tau E^{}(\overline{\psi }).`$ $`E^{}`$ needs to satisfy the conditions:
$`a)`$ $`E^{}(\overline{\psi }e_1)=\tau E^{}(\overline{\psi })`$ (47)
$`b)`$ $`E^{}(\overline{\psi }e_i)=E^{}(\overline{\psi })e_ii>2`$ (48)
$`c)`$ $`E^{}(\overline{\psi }e_1)=0\overline{\psi }e_1=0.`$ (49)
ยฟFrom (46), in order to identify $`\overline{\sigma }`$ with $`\sigma `$, we see that the projection $`E^{}`$ must satisfy:
$`E^{}(F_{\omega _n}\overline{\sigma })=F_{\omega _{n2}}.`$ (50)
$`E^{}`$ can be realized as:
$`E^{}(\overline{\psi })=c^{}{\displaystyle \frac{1}{\varphi (z,z_i)}}\overline{\psi }(z_1=z,z_2=q^2z,z_i),`$ (51)
where $`\varphi (z,z_i)=_{i=3}^n(z_iq^4z)`$ and $`c^{}`$ is a normalization constant. $`E^{}`$ verifies (49a,b) by construction as can be seen from the expression (42) of $`e_1\tau `$.
Using the explicit expression (27) of $`\sigma `$, we have:
$`E^{}(F_{\omega _n}\sigma )=c^{}s^{\frac{n}{4}\frac{1}{2}}{\displaystyle \frac{1}{\varphi (z,z_i)}}{\displaystyle \underset{3}{\overset{\frac{n}{2}+1}{}}}(q^3zq^1z_i){\displaystyle \underset{\frac{n}{2}+2}{\overset{n}{}}}(qz_iq^1sz)F_{\omega _{n2}}(z_3,\mathrm{},z_n).`$ (52)
which imposes $`s=q^6`$ for the polynomial in the numerator to be proportional to $`\varphi (z,z_i)`$ and (50) to be satisfied.
To identify $`_n`$ and $`_n^{}`$, we give a more convenient characterization of $`_n^{}`$. Consider the space $`_n^{\prime \prime }`$ of homogenous polynomials in $`n`$ variables, and of the minimal total degree, obeying the property:
$`(\mathrm{P}):\overline{\psi }(z_i=z,z_j=q^2z,z_k=q^4z)=0,\mathrm{if}i,j,k,\mathrm{are}\mathrm{cyclically}\mathrm{ordered}.`$ (53)
This property is obviously compatible with the cyclic identification $`z_{i+n}=q^6z_i`$, it is thus preserved by $`\overline{\sigma }`$ (27). By applying (P) to the triplets $`(1,2,j)`$, we see that the projection (51) is well defined from $`_n^{\prime \prime }`$ to $`_{n2}^{\prime \prime }`$.
We show that $`_n^{\prime \prime }=_n`$. For this, we first show that $`_n^{\prime \prime }`$ is a module over the T.L. algebra $`๐_n`$ and that it contains $`_n`$, then we show that $`_n^{\prime \prime }`$ is irreducible over $`๐_n`$.
To show that $`_n^{\prime \prime }`$ is a module over $`๐_n`$, we verify that the generators $`e_i`$ preserve the property (P). Assuming that the polynomial $`\overline{\psi }`$ verifies (P) we verify that $`\overline{\psi }e_i`$ obeys (P) for a cyclically ordered triplet $`k,l,m`$. If $`\{i,i+1\}\{k,l,m\}=\mathrm{}`$, it is obvious. If $`i+1=k`$, it results from the fact that $`\overline{\psi }`$ obeys (P) for the triplets $`i,l,m`$ and $`i+1,l,m`$. The same type of argument applies if $`i=m`$. If $`\{i,i+1\}\{k,l,m\}`$, $`\overline{\psi }(e_i\tau )`$ is proportional to $`(qz_iq^1z_{i+1})`$ and therefore obeys (P).
Let us show that (49c) is satisfied in $`_n^{\prime \prime }`$. If $`E^{}(\overline{\psi }e_1)=0`$, $`\overline{\psi }e_1`$ vanishes when $`z_2=q^2z_1`$, and from the definition (42) of $`e_1`$, it is symmetric in $`z_1,z_2`$. It is therefore divisible by $`(z_1q^2z_2)(z_2q^2z_1)`$. Hence, $`\overline{\psi }e_1/(z_1q^2z_2)`$ satisfies (P) and has a total degree reduced by one. It is thus equal to zero according to our minimal degree hypothesis.
It is clear that $`F_\omega `$ satisfies the property (P). To show that $`_n_n^{\prime \prime }`$, we need to show that the degree of the polynomials in $`_n^{\prime \prime }`$ is the degree $`\frac{n}{2}(\frac{n}{2}1)`$ of $`F_\omega `$. We proceed by recursion on $`n`$ and for the moment, we exclude the case where $`e_1`$ is represented as zero in $`_n^{\prime \prime }`$. Due to (49c) there are polynomials $`\overline{\psi }`$ in $`_n^{\prime \prime }`$ such that $`E^{}(\overline{\psi })0`$. This implies that $`\overline{\psi }`$ has a degree at least $`n2`$ in $`z_1,z_2`$. We can apply the recursion hypothesis to $`E^{}(\overline{\psi })_{n2}^{\prime \prime }`$ to conclude that the minimal degree is $`\frac{n}{2}(\frac{n}{2}1)`$.<sup>2</sup><sup>2</sup>2The same argument shows that the maximal degree of the polynomials in $`_n^{\prime \prime }`$ is $`\lambda `$ for the order defined in D.3.
To show that $`_n^{\prime \prime }`$ is irreducible as a T.L. module, we use the recursion hypothesis that $`_{n2}^{\prime \prime }=_{n2}`$. Due to (49c), $`E^{}`$ is injective from $`_n^{\prime \prime }e_1`$ to $`E^{}(_n^{\prime \prime })_{n2}`$. Since $`_ne_1=_{n2}_n^{\prime \prime }e_1`$, we have $`_n^{\prime \prime }e_1=_ne_1`$. Thus, if $`_n^{\prime \prime }`$ contains an irreducible submodule $`R_n`$, $`Re_1=0`$. If $`Re_1=0,`$ from (8) we see that all the $`e_i`$ are represented as $`0`$ in $`R`$, and therefore, the polynomials in $`R`$ are proportional to $`\mathrm{\Delta }_n`$ defined in (43) times a symmetric polynomial. Since the total degree of $`\mathrm{\Delta }_n`$ is larger than $`\frac{n}{2}(\frac{n}{2}1)`$, $`R=0`$. We conclude that $`_n^{\prime \prime }=_n`$ as a T.L. module.
Finally, to identify $`_n^{\prime \prime }`$ and $`_n`$ as affine modules, we observe that $`y_1=\sigma ^1\overline{\sigma }`$ commutes with $`๐_{n1}`$ generated by $`e_2,\mathrm{},e_n`$. Since $`_n`$ is irreducible over $`๐_{n1}`$ , $`y_1`$ is proportional to the identity, thus $`\sigma `$ and $`\overline{\sigma }`$ can be identified.
#### 4.3.1 Relation with the Macdonald Polynomials and the work of Di Francesco and Zinn-Justin.
As a check of consistency, we must verify that the two expressions of the eigenvalue of the central operators $`y`$ (30,35) are the same when $`s=q^6`$. This is indeed the case if we substitute in (30) the degree $`\lambda `$ (31) of the highest polynomial in $`_n`$ and $`c=q^{3(1\frac{n}{2})}`$.
For a generic $`s`$, the operator $`y`$ (30) can be diagonalized on the basis of symmetric polynomials and its eigenvectors define the Macdonald polynomials . We have seen that when $`s=q^6`$, the polynomial representation is reducible. As a counterpart, some diagonal elements $`y_\lambda ^{}`$ of $`y`$ become degenerate with $`y_\lambda `$, for example, $`\lambda _2^{}=\lambda _21`$, $`\lambda _5^{}=\lambda _5+1`$. Thus, $`y`$ cannot be diagonalized. We must use another operator such as $`\frac{dy}{ds}`$ to define the analogous symmetric polynomial.
In the non semisimple case $`q^2+q+1=0,(\tau =1),`$ the T.L. representation admits a sub-representation given by $`\mathrm{Vec}\{x_\pi \pi ,\mathrm{with}x_\pi =0\}.`$ The trivial representation $`\mathrm{\Omega }`$ is obtained by equating to zero these vectors. The dual polynomial $`F_\mathrm{\Omega }=_\pi F_\pi `$ is therefore symmetrical of degree $`\lambda `$, and obeys the property (P)(53). This completely determines it to be proportional to the Schur function $`s_\lambda `$ with $`\lambda `$ given by (31). Indeed, $`s_\lambda `$ has a degree $`\lambda `$ and satisfies (P) since three columns of the determinant which defines it become linearly dependant when we make the substitution (P). By the same argument as used in 4.3, the degree of a symmetric polynomial satisfying (P) must be at least $`\lambda `$ (relatively to the order of partitions which follows from D.3), and this proves its uniqueness.
In this limit, the $`F_\pi `$ are also the components the ground state of the $`O(n=1)`$ transfer matrix, and thus, we have proved that this sum is $`s_\lambda `$, which is the result of .
It would be interesting to see if in this limit, $`F_\mathrm{\Omega }`$ can be recovered as the eigenvector of some operator such as $`\frac{dy}{ds}`$.
### 4.4 Representation on Gaudinโs determinants.
It is well known that the Bethe scalar products can be expressed using a quotient of two determinants. Here, we construct a representation of the T.L. algebra acting on these quotients. We split the variables $`z_i`$ into $`A=\{z_1,\mathrm{},z_{\frac{n}{2}}\}`$ and $`B=\{z_{\frac{n}{2}+1},\mathrm{},z_n\}`$. We also introduce $`p`$ a square root of $`q`$, $`p^2=q`$. We define the polynomial $`F_\omega ^{}`$:
$`F_\omega ^{}(z_1,..,z_n)={\displaystyle \frac{\left|(p^2z_ip^2z_j)^1(pz_ip^1z_j)^1\right|}{\left|(pz_ip^1z_j)^1\right|}}\mathrm{\Delta }_n^q(z_1,\mathrm{},z_n),\mathrm{with}iA,jB.`$ (54)
The first factor is the ratio of the Gaudin determinant with the Cauchy determinant . It is also related to the domain wall boundary condition partition function of a six vertex model with weights: $`a=qxq^1y,b=pxp^1y,c=\sqrt{xy}(pp^1)`$ <sup>3</sup><sup>3</sup>3Notice that for this six vertex model, $`\frac{a^2+b^2c^2}{ab}=p+p^1\tau `$..
The second factor $`\mathrm{\Delta }_n`$ (43) transforms the scalar product into a polynomial. This factor has an innocuous effect on the T.L. algebra since:
$`\mathrm{\Delta }_n^q(z_1,\mathrm{},z_n)t_i=\stackrel{~}{t}_i\mathrm{\Delta }_n^q(z_1,\mathrm{},z_n),`$ (55)
where $`\stackrel{~}{t}_i`$ is obtained from $`t_i`$ (28) by the substitution $`q1/q`$ which preserves the relations (3,4), but exchanges $`e_i`$ with $`e_i\tau `$.
The ratio of the two determinants being symmetrical in the two sets of variables $`A`$ and $`B`$, $`F_\omega ^{}`$ is annihilated by all the $`e_i`$ with $`i\frac{n}{2}`$.
To show that the action of the T.L. algebra (42) on $`F_\omega ^{}`$ produces an irreducible module, we proceed as in 4.3. Consider the space $`_n`$ of homogenous polynomials in $`n`$ variables, and of the minimal total degree, obeying the property:
$`(\mathrm{P}^{}):`$ $`\overline{\psi }(z_{i_1}=q^{a_1}z,z_{i_2}=q^{a_2}z,z_{i_3}=q^{a_3}z)=0,\mathrm{if}i_1,i_2,i_3,\mathrm{are}\mathrm{cyclically}\mathrm{ordered},`$ (57)
$`\mathrm{and}\mathrm{for}:(a_1,a_2,a_3)=(1,0,1),(1,1,0),(2,0,2),(0,1,1).`$
Note that these triplets are stable under the cyclic permutation, $`(a_1,a_2,a_3)(a_32,a_2+1,a_1+1)`$, and the transpositions, $`(a_i,a_{i+1})(a_{i+1},a_i),`$ whenever $`|a_{i+1}a_i|=1`$.
ยฟFrom the cyclic invariance, we deduce that this space is preserved under the action of $`\sigma `$ (27) if we take $`s=q^3`$.
By applying the property (Pโ) to $`z_1,z_2,z_i`$ with $`(a_1,a_2,a_3)=(1,1,0)`$ and $`(2,0,2)`$, we can define a projection (51) from $`_n_{n2}`$. The polynomial $`\varphi (z,z_i)`$ is now a product of two factors $`\varphi (z,z_i)=_{i=3}^n(qzz_i)(q^4zz_i)`$. Arguing as in 4.3, we see that this projection satisfies the properties (49).
This space is stable under the action of the generators $`e_i`$. The proof is similar to the one given in 4.3 and requires the stability of the triplets $`(a_1,a_2,a_3)`$ under the transpositions. The minimal degree is now $`n(\frac{n}{2}1)=2|\lambda |`$ with $`|\lambda |`$ given by (31).
Let us show that $`F_\omega ^{}`$ (54) satisfies the property (Pโ) (57). We consider $`(i_1,i_2,i_3)`$ and $`(a_1,a_2,a_3)`$. If the variables $`z_l,z_m`$ with $`l<m`$, corresponding to two $`a_i`$ which differ by $`2`$, belong to the same set $`A`$ or $`B`$, $`F_\omega ^{}(z_m=q^2z_l)=0`$ due to the factor $`\mathrm{\Delta }_n`$. Otherwise, two variables $`z_l=zA`$ and $`z_m=q^2zB`$ differ by a factor $`q^2`$. By isolating the contribution of the pole $`(p^2z_lp^2z_m)`$ in the Gaudin determinant, we factorize a term $`_i(qzz_i)`$ coming from the Cauchy denominator, and this enables to conclude that $`F_\omega ^{}(z_{i_1}=q^{a_1}z,z_{i_2}=q^{a_2}z,z_{i_3}=q^{a_3}z)=0`$ in all the other cases.
Arguing as in 4.3 we conclude that $`_n`$ is an irreducible module over the affine T.L. algebra and that it coincides with the module obtained upon acting with the generators on $`F_\omega ^{}`$.
We verify again that the eigenvalue of the central operators $`y`$ (30) is given by (35). Now, $`s=q^3`$ instead of $`q^6`$ in 4.3, but the degree $`2\lambda `$ (31) of the highest polynomial in $`_n`$ is doubled and $`c`$ keeps the same value $`c=q^{3(1\frac{n}{2})}`$.
In the nonsemisimple case $`q^2+q+1=0`$, using the result of we have $`F_\pi ^{}=s_\lambda F_\pi `$, and therefore, $`_\pi F_\pi ^{}=s_\lambda ^2`$.
## 5 Conclusion.
Let us conclude with a few comments and questions.
On the mathematical side, this work provides a unification ground around the conjectures relating the eigenvector components of a loop model transfer matrix, the six vertex model domain wall boundary condition partition function and other mathematical objects. It opens the possibility to deform the polynomials underlying these conjectures by presenting them from the algebra representation point of view. We believe that these conjectures are related to incompressibility, and we hope to return to this point in a future publication.
ยฟFrom a technical point of view, it would be interesting to repeat the Jones construction of A.2 on the polynomials directly. This would allow to recover in a direct way the product structure which they carry since they are associated to elements of the T.L. algebra.
The precise correspondence between the polynomial obtained here and the Macdonald polynomials needs to be clarified.
Finally, do the deformed wave functions considered here have anything to do with physics? At this moment, we have no answer to this question. A step towards a physical interpretation would be to identify a scalar product and a Hermitian Hamiltonian for which these wave functions are the ground states. This could also be useful to access to the excited states which play an important role in the Quantum Hall Effect.
### 5.1 Acknowledgements
I wish to thank Philippe di Francesco for generously explaining me his works and for discussions.
I am greatly indebted to Kirone Mallick, Gregoire Misguich and particularly Bertrand Duplantier for their help during the course of this work.
## Appendix A Word representation.
### A.1 Reduced words.
The module $`_n`$ is obtained by acting with the T.L. generators of $`๐_n`$ on the lowest state $`\alpha =e_1e_3\mathrm{}e_{n1}`$. Using the relations (8), we obtain a basis of $`_n`$ given by reduced words $`\pi `$:
$`\pi =(e_{a_{n1}}e_{a_{n1}+1}..e_{n1})\mathrm{}(e_{a_{2p+1}}e_{a_{2p+1}+1}..e_{2p+1})\mathrm{}(e_{a_3}e_{a_3+1}..e_3)e_1,`$ (58)
with, $`a_{2p+1}2p+1`$, and $`1<a_3<..<a_{2p+1}<..<a_{n1}`$. So, a word is fully characterized by the sequence $`(a_{2p+1})`$.
On reduced words there is a natural order relation: $`\pi >\pi ^{}`$ if $`\pi `$ is written $`b\pi ^{}`$ with b a monomial. One has $`\pi \pi ^{}`$ if $`a_{2p+1}a_{2p+1}^{}`$ for all $`p`$.
Another way to represent a reduced word is in terms of paths. Let $`m_i`$ be the number of times the generator $`e_i`$ appears in the reduced expression of $`\pi `$. One has $`m_{2i}=m_{2i1}\mathrm{or}m_{2i1}1`$ and $`m_{2i+1}=m_{2i}\mathrm{or}m_{2i}+1`$. We define $`h_{2i}=2m_{2i}1`$, $`h_{2i1}=2m_{2i1}2`$ and $`h_0=h_n=0`$ by convention. We can describe the words $`\pi `$ by the paths $`\pi =[h_i]`$ where $`h_0=h_n=0`$, $`h_i0`$ and $`h_{i+1}h_i=\pm 1`$. Using the path representation, one has $`\pi \pi ^{}`$, if $`[h_i][h_i^{}]i.`$
If $`\pi ^{}`$ is a non reduced word, by reducing it, one decreases the number of times the generator $`e_i`$ appears in its expression. We thus see that the order relation can be presented in a weaker form valid for non reduced words: If $`\pi ^{}`$ is a word, not necessarily reduced and $`\pi `$ is a reduced word, $`\pi >\pi ^{}`$ if $`\pi ^{}`$ can be obtained by erasing letters $`e_k`$ from the (reduced) expression of $`\pi `$.
Finally, there is way to characterize this representation in terms of link patterns. It is convenient to dispose the $`n`$ points cyclically around a circle. A link pattern is obtained by pairing all the points in the set $`\{1,2,\mathrm{},n\}`$: $`\pi =\{[i_1,i_2],[i_3,i_4],\mathrm{},[i_{n1},i_n]\}`$, in such a way that two links never cross. In practise, if $`[i,j]`$ is a link, then the other links $`[k,l]`$ are either inside, or outside the interval $`[i,j]`$. The action of $`e_{i,i+1}`$ is given by: $`e_{i,i+1}\{[i,i+1],\mathrm{},[i_{n1},i_n]\}=\tau \{[i,i+1],\mathrm{},[i_{n1},i_n]\}`$, and $`e_{i,i+1}\{[i,j],[i+1,k]\mathrm{},[i_{n1},i_n]\}=\{[i,i+1],[j,k],\mathrm{}\}`$. In this representation, $`\alpha =\{[1,2],[3,4],\mathrm{},[n1,n]\}`$, and $`\omega =\{[1,n],[2,n1],\mathrm{},[\frac{n}{2}1,\frac{n}{2}+1]\}`$.
These representations are illustrated in figure 1.
### A.2 Identifying $`_n`$ with $`๐_{\frac{n}{2}}.`$
The link pattern representation allows to identify in a natural way $`_n`$ with $`๐_{\frac{n}{2}}.`$ If we split $`\{1,2,\mathrm{},n\}`$ into two subsets :$`\{1,2,\mathrm{},\frac{n}{2}\}`$ and $`\{\frac{n}{2},\mathrm{},n\}`$, the product $`\pi \pi ^{}`$ is defined on the link patterns by โstackingโ the two link patterns and identifying the last $`\frac{n}{2}`$ points of $`\pi `$ with the first $`\frac{n}{2}`$ points of $`\pi ^{}`$ through $`in+1i`$. The link pattern $`\pi \pi ^{}`$ is obtained by removing the loops which appear in this concatenating operation by giving them a weight $`\tau `$.
Another identification can be achieved on paths by folding a path of length $`n`$ into a loop of length $`\frac{n}{2}`$. In this way, we realize $`๐_{\frac{n}{2}}`$ as the algebra of double paths acting on Bratteli diagrams .
In this identification, $`๐_{\frac{n}{2}}`$ is a bimodule over itself. The first $`\frac{n}{2}1`$ generators $`e_i๐_n`$ are identified with the generators of $`๐_{\frac{n}{2}}`$ acting to the left, while the last $`\frac{n}{2}1`$ generators are identified with $`e_{n+1i}`$ acting to the right.
The state $`\omega `$ is the identity in $`๐_{\frac{n}{2}}`$, and the trace in $`๐_{\frac{n}{2}}`$ coincides with the scalar product with $`\omega `$ in $`๐_n`$:
$`\mathrm{tr}(x)=\tau ^{\frac{n}{2}}\omega |x.`$ (59)
The projection: $`E_{\frac{n}{2}}=\sigma ^{\frac{n}{2}+1}E\sigma ^{\frac{n}{2}1},`$ with $`E`$ given by (36) can be reinterpreted as a conditional expectation value , $`E_{\frac{n}{2}}:๐_{\frac{n}{2}}๐_{\frac{n}{2}1}`$. Jones construction enables then to construct $`e_{\frac{n}{2}}๐_{\frac{n}{2}+1}`$ algebraically from the knowledge of $`E_{\frac{n}{2}}`$.
### A.3 Triangularity of $`y_m`$.
Let us show that the affine generators $`y_{m+1}=t_m^1t_{m1}^1\mathrm{}t_1^2\mathrm{}t_m^1`$ are triangular in the word representation. It is obvious for $`y_1=1`$ and $`y_2=t_1^2`$ since $`e_1`$ is triangular. We proceed by recursion and assume that $`y_k`$ are triangular for $`k<m+1`$. Using these hypotheses, we show that $`y_{m+1}`$ is also triangular.
First we show that $`y_{m+1}`$ acts diagonally on $`\alpha `$. To study the action of $`y_{m+1}`$ on $`\alpha `$, we distinguish the two cases $`m`$ odd or even. If m is odd, then:
$`y_{m+1}\alpha =t_m^1y_mt_m^1e_m\mathrm{}=qt_m^1y_me_m\mathrm{}=\lambda _mqt_m^1e_m\mathrm{}=q^2\lambda _m\alpha ,`$ (60)
where $`\lambda _m`$ is the eigenvalue of $`y_m`$ on $`\alpha `$. If $`m`$ is even, we make use of the fact that $`t_{m1}^1t_m^1e_{m1}=\frac{1}{q}e_me_{m1}`$ and the same relation with the indices $`m`$ and $`m1`$ exchanged to obtain:
$`y_{m+1}\alpha `$ $`=`$ $`t_m^1t_{m1}^1y_{m1}t_{m1}^1t_m^1e_m\mathrm{}={\displaystyle \frac{1}{q}}t_m^1y_{m1}e_me_{m1}\mathrm{}`$ (61)
$`=`$ $`{\displaystyle \frac{1}{q}}\lambda _{m1}t_m^1t_{m1}^1e_me_{m1}\mathrm{}={\displaystyle \frac{1}{q^2}}\lambda _{m1}\alpha .`$ (62)
We deduce that $`\alpha `$ is an eigenstate of $`y_m`$ with the eigenvalue $`\lambda _m`$ obeying the recursion relations $`\lambda _{2m}=q^2\lambda _{2m1}`$, $`\lambda _{2m+1}=\frac{1}{q^2}\lambda _{2m1}`$. Together with the fact that $`\lambda _1=1`$, we deduce (30).
To show that $`y_{m+1}`$ is triangular on words $`\alpha `$. We proceed by recursion and assume that $`y_{m+1}`$ acts in a triangular way on words $`<\pi `$ and show that the property is also true for $`\pi `$.
Let us consider the action of $`y_{m+1}`$ on a reduced word $`\pi \alpha `$. This word can be put under the form $`\pi =e_i\pi ^{}`$ where $`\pi ^{}<\pi `$. We consider the three cases, $`im,m+1`$, $`i=m`$, $`i=m+1`$. In the third case, either the word can be written in the form $`e_{m+1}e_m\pi ^{}`$ with $`\pi ^{}`$ reduced, or it can be written $`e_p\pi ^{}`$ with $`p<m`$. The second possibility reduces to the first case and we need only consider the first possibility.
We observe that $`y_{m+1}`$ commutes with $`e_i`$: $`y_{m+1}e_i=e_iy_{m+1}`$ if $`i>m+1`$ or if $`i<m`$. It is obvious if $`i>m+1`$ and follows from the braid relations if $`i<m`$. In the three cases we can thus write:
$`y_{m+1}e_i\pi ^{}`$ $`=`$ $`e_i(y_{m+1}\pi ^{})\mathrm{for}im,m+1,`$ (63)
$`y_{m+1}e_m\pi ^{}`$ $`=`$ $`t_m^1y_mt_m^1e_m\pi ^{}=qt_m^1(y_me_m\pi ^{}),`$ (64)
$`y_{m+1}e_{m+1}e_m\pi ^{}`$ $`=`$ $`t_m^1y_mt_m^1e_{m+1}e_m\pi ^{}=t_m^1(y_me_m\pi ^{}+{\displaystyle \frac{1}{q}}y_me_{m+1}e_m\pi ^{}).`$ (65)
It follows from the hypothesis that the terms in brackets are less than $`\pi `$. In the first case because $`\pi ^{}<\pi `$, and in the two others because $`y_m`$ is assumed to be triangular.
To conclude that $`y_{m+1}`$ is triangular, we must show that the action of $`e_i`$ in the first case and $`e_m`$ in the two other cases preserves the triangularity : If $`e_i\pi `$ is a reduced word and $`\pi ^{}\pi `$, then, $`e_i\pi ^{}e_i\pi `$. If $`e_m\pi `$ is a reduced word and $`\pi ^{}e_m\pi `$, then $`e_m\pi ^{}e_m\pi `$. Finally, if $`e_{m+1}e_m\pi `$ is a reduced word and $`\pi ^{}e_{m+1}e_m\pi `$, then $`e_m\pi ^{}e_{m+1}e_m\pi `$. These properties follow from the weak form of the order relation. This concludes the proof of triangularity of $`y_{m+1}`$.
### A.4 Action of $`\sigma `$ on words.
The action of $`\sigma =q^{\frac{n}{2}2}t_{n1}^1\mathrm{}t_1^1`$ on words can be computed similarly. First, using the braid relation (3), one sees that $`\sigma e_i=e_{i1}\sigma `$ for $`i>1`$. To fully characterize its action, we must compute $`(\sigma \alpha )`$. Using $`t_1^1e_1=qe_1`$ and $`t_{m+1}^1t_m^1e_{m+1}=\frac{1}{q}e_me_{m+1}`$ we obtain:
$`\sigma \alpha ={\displaystyle \underset{i=1}{\overset{\frac{n}{2}1}{}}}e_{2i}\alpha .`$ (66)
Thus, $`(\sigma \alpha )`$ can be characterized by the property:
$`e_{2i}(\sigma \alpha )=\tau (\sigma \alpha ),`$ (67)
for $`1i\frac{n}{2}`$. $`(\sigma \alpha )`$ can then be used as a lowest state to construct a canonical basis by acting on it with the generators $`e_2,\mathrm{},e_n`$.
## Appendix B Explicit construction of $`_4`$.
Let us construct $`_4`$ the dual of $`_4`$. The basis of $`_4`$ is given by the words $`e_1e_3=\alpha ,e_2e_1e_3=\omega `$. So we search for a vector $`\mathrm{\Psi }`$ of the form:
$`\mathrm{\Psi }=F_\alpha (z_1,..,z_4)\alpha +F_\omega (z_1,..,z_4)\omega ,`$ (68)
where $`F_\alpha ,F_\omega `$ are polynomials of degree $`(1,1)`$ in the variables $`z_i`$. The action of the $`T.L.`$ affine algebra is given by the matrices:
$`e_1=e_3=\left(\begin{array}{cc}\tau & 1\\ 0& 0\end{array}\right),e_2=e_4=\left(\begin{array}{cc}0& 0\\ 1& \tau \end{array}\right),\sigma =\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right).`$ (69)
We can obtain the dual representation by acting with the generators on $`F_\omega \left(\begin{array}{c}0,1\end{array}\right)`$ annihilated by $`\overline{e}_1=e_1,\overline{e}_3=e_3`$. The minimum degree polynomial annihilated by $`e_1,e_3`$ is given by:
$`F_\omega =(qz_1q^1z_2)(qz_3q^1z_4).`$ (70)
Let us take $`\overline{\sigma }=\sigma `$ of the form:
$`F(z_1,z_2,z_3,z_4)\sigma =cF(z_2,z_3,z_4,sz_1).`$ (71)
We obtain two different expression for $`F_\alpha (1,0)`$ which we must equate. One results from the cyclic property: $`F_\alpha =F_\omega \sigma `$, the other given by: $`F_\alpha =F_\omega (e_2\tau )`$.
We get the equation:
$`{\displaystyle \frac{(qz_1q^1z_2)(qz_3q^1z_4)(qz_1q^1z_3)(qz_2q^1z_4)}{z_2z_3}}=c(qz_4q^1sz_1),`$ (72)
which determines $`s=q^6`$, $`c=q^3`$, and:
$`F_\alpha =(qz_2q^1z_3)(q^2z_4q^2z_1).`$ (73)
## Appendix C Module $`F_\omega `$.
Let us define a T.L. module $`M`$ defined in terms of a state $`F_\omega `$ satisfying $`F_\omega e_i=0`$ for $`i\frac{n}{2}`$. The module is obtained by acting with the T.L. generators and reducing words using the T.L. relations (8). In this module, a canonical basis is:
$`\overline{\psi }=F_\omega (e_{\frac{n}{2}}e_{\frac{n}{2}1}\mathrm{}e_{a_{\frac{n}{2}}+1}e_{a_{\frac{n}{2}}})\mathrm{}(e_pe_{p1}\mathrm{}e_{a_p})\mathrm{}(e_{n1}\mathrm{}e_{a_{n1}}),`$ (74)
where the $`p`$ take the all the values between $`\frac{n}{2}`$ and $`n1`$ and the $`a_p`$ are restricted by the conditions: $`a_pp+1`$, $`a_{\frac{n}{2}}<a_{\frac{n}{2}+1}..<a_p<..<a_{n1}`$. The convention is that if $`a_p=p+1`$, the sequence $`(e_p\mathrm{}e_{a_p})`$ is empty. A word $`\overline{\psi }`$ is fully characterized by the sequence $`(a_p)`$. The word can also be associated to the Young diagram $`[\mu _{p+1\frac{n}{2}}]=[pa_p+1]`$.
There is a reflection symmetry, $`ini`$, and an alternative description of the module in terms of reflected words:
$`\overline{\psi }=F_\omega (e_{\frac{n}{2}}\mathrm{}e_{b_{\frac{n}{2}}1}e_{b_{\frac{n}{2}}})\mathrm{}(e_p\mathrm{}e_{b_p1}e_{b_p})\mathrm{}(e_1\mathrm{}e_{b_1}),`$ (75)
$`1p\frac{n}{2}`$, $`b_pp1`$, $`b_{\frac{n}{2}}>\mathrm{}>b_1`$. It is associated to the dual Young diagram $`[\mu _{\frac{n}{2}p+1}^{}]=[b_pp+1]`$
A similar order relation as defined earlier holds for reduced words, $`\overline{\psi }^{}<\overline{\psi }`$ if $`\overline{\psi }`$ can be written $`\overline{\psi }=\overline{\psi }^{}a`$. For non reduced words $`\overline{\psi }^{}`$, it is sufficient that $`\overline{\psi }^{}`$ can be obtained by erasing letters $`e_k`$ from the (reduced) expression of $`\overline{\psi }`$.
In general, the module $`F_\omega `$ is reducible, it will be irreducible if $`F_\omega `$ satisfies the Fock condition:
$`F_\omega (1+{\displaystyle \underset{m=0}{\overset{\frac{n}{2}1}{}}}q^{m+1}t_{\frac{n}{2}}\mathrm{}t_{\frac{n}{2}m})=0.`$ (76)
In this case, The only allowed words $`\overline{\psi }`$ (74) can be associated to their complementary $`\pi _\psi `$ in such a way that one can write without reducing the expression:
$`\psi \pi _\psi =\omega .`$ (77)
Thus, we get the supplementary constraint $`a_p>2p+1n,b_p<2p1`$.
### C.1 Reducing the Hecke Module to its T.L. form.
Let us consider a module $`M^{}`$ over the Hecke algebra defined by acting with the Hecke algebra generators satisfying (7) on the state $`F_\omega `$ satisfying $`F_\omega e_i=0`$ for $`i\frac{n}{2}`$ . We want to show that the Hecke algebra acts as a T.L. algebra on this module. For this, we first show that the Hecke relations (7) are sufficient to reduce the word basis of $`M^{}`$ to be of the T.L. form (74). Thus, $`M^{}`$ and $`M`$ can be identified as vector spaces. From this, we will deduce that $`M^{}=M`$ as modules. In other words, the projectors $`U_{i,i+1}^{}=e_ie_{i+1}e_ie_i`$ are null in $`M^{}`$.
Let us assume that it is not true. Since all the basis elements of $`M^{}`$ are obtained upon acting on $`F_\omega `$ with letters $`e_k`$, there is a basis element $`\overline{\psi }e_i`$ which cannot be expressed as a linear combination of words of the form (74) although $`\overline{\psi }`$ is of the form (74). Among all the $`\overline{\psi }`$ which verify this property, we can take the smallest possible for the order relation, so that that $`\overline{\psi }^{}e_i`$ is of the form (74) when $`\overline{\psi }^{}<\overline{\psi }`$. We show that this leads to a contradiction.
Let us consider the word $`\overline{\psi }e_i`$. It is a word of the form (74) in the three following cases. When $`\overline{\psi }e_i`$ is a reduced word $`>\overline{\psi }`$, for $`i=a_p1`$ if $`a_p1>a_{p1}`$. When $`\overline{\psi }e_i=\tau \overline{\psi }`$ when $`i=a_p`$ and $`a_p>a_{p+1}1`$. When $`\overline{\psi }e_i=0`$ if $`i<a_{\frac{n}{2}}1`$ or $`i>b_{\frac{n}{2}}+1`$.
The two remaining cases to consider are: First, when $`a_p<i<a_{p+1}1`$ for some $`p`$. Second, when $`a_p<ia_p+k`$ if $`a_{p+k}=a_p+k`$ with $`k1`$. The second case can be studied similarly to the first one using the reflection symmetry $`ini`$ and corresponds to $`b_p^{}>i>b_{p^{}1}+1`$.
In the first case, $`\overline{\psi }e_i=\overline{\psi }^{}(e_p\mathrm{}e_{a_p+1}e_{a_p})e_i(e_{p+1}\mathrm{}e_{a_{p+1}})\mathrm{}`$, and using the relation (7), we see that:
$`e_p\mathrm{}e_{a_p+1}e_{a_p}e_i=e_{i1}e_p\mathrm{}e_{a_p+1}e_{a_p}+e_p\mathrm{}e_{i+1}(e_ie_{i1})e_{i2}\mathrm{}e_{a_p+1}e_{a_p}.`$ (78)
The second term is $`<\overline{\psi }`$ and therefore of the T.L. form by the recursion hypothesis. The first term can be eliminated by repeating this relation $`p\frac{n}{2}`$ times to push $`e_i`$ and then $`e_{i1},\mathrm{},e_{i+\frac{n}{2}p}`$ to the left of the word. The last application of the relation gives a term $`F_\omega e_{i+\frac{n}{2}p1}=0`$ since $`i+\frac{n}{2}p1<\frac{n}{2}`$.
This exhaust all the possibilities and $`\overline{\psi }e_i`$ can always be expressed as a linear combination of reduced T.L. words (74) in contradiction with the hypothesis. Therefore, the word basis of $`M^{}`$ coincides with the word basis (74).
To conclude that $`M^{}=M`$, let us consider the projectors $`U_{i,i+1}^{}=e_ie_{i+1}e_ie_i`$, and the space $`M^{\prime \prime }M^{}`$ annihilated by all the $`U_{i,i+1}^{}`$. The space $`M^{\prime \prime }`$ defines a module for the T.L. algebra. Since $`F_\omega M^{\prime \prime }`$, this module can be identified with $`M`$. Therefore, $`M`$ is a subspace of $`M^{}`$ with the same dimension, and thus, $`M=M^{}`$.
## Appendix D Yang-Baxter Equation and Polynomials:
### D.1 Polynomial representation of the Hecke generators
In this section, we derive the expression of the Hecke generators $`\overline{t}_i`$ (28) from the Yang-Baxter equation.
The Yang-Baxter algebra (also called $`RLL=LLR`$ relation) can be expressed as:
$`R_{12}(z_1,z_2)L_1(z_1)L_2(z_2)=L_2(z_2)L_1(z_1)R_{12}(z_1,z_2),`$ (79)
where $`R_{12}(z_1,z_2)`$ is a solution of the Yang-Baxter equation:
$`R_{12}(z_1,z_2)R_{13}(z_1,z_3)R_{23}(z_2,z_3)=R_{23}(z_2,z_3)R_{13}(z_1,z_3)R_{12}(z_1,z_2).`$ (80)
If we assume that $`R_{12}(z_1,z_2)=Y_{12}(z_1,z_2)P_{12}`$ where $`P_{12}`$ acts in the natural way on the spin indices, $`(P_{12}t_{13}=t_{23}P_{12}),`$ but commutes with $`z_i`$, (79) rewrites as:
$`Y_{12}(z_1,z_2)L_2(z_1)L_1(z_2)=L_2(z_2)L_1(z_1)Y_{12}(z_1,z_2)=L_2(z_1)L_1(z_2)k_{12},`$ (81)
where $`k_{12}`$ acts to the left by permuting the variables $`z_1,z_2`$. If we normalize of $`Y(z_1,z_2)`$ so that:
$`Y_{12}(z_1,z_2)Y_{12}(z_2,z_1)=1,`$ (82)
it is consistent to demand that the $`Y_{ii+1}`$ act as a representation of the permutation algebra on some wave function $`\mathrm{\Psi }`$:
$`Y_{12}(z_1,z_2)\mathrm{\Psi }(z_1,z_2)=\mathrm{\Psi }(z_2,z_1)=\mathrm{\Psi }(z_1,z_2)k_{12}.`$ (83)
The $`Y_{ij}`$ are are called Yangโs operators in .
A well known solution of (80),(82) in terms of the Hecke algebra (4) is:
$`Y_{12}(z)={\displaystyle \frac{t_{12}zt_{12}^1}{zqq^1}},`$ (84)
where $`z=\frac{z_1}{z_2}`$.
Substituting (84) in (83), we can also rewrite this relation as:
$`t_{12}\mathrm{\Psi }(z_1,z_2)=\mathrm{\Psi }(z_1,z_2)\overline{t}_{12},`$ (85)
Where $`\overline{t}_{12}`$ takes the form:
$`\overline{t}_{12}=q^1+(1k_{12}){\displaystyle \frac{qz_1q^1z_2}{z_1z_2}},`$ (86)
and this coincides with (28).
### D.2 Commutation relations of the affine generators $`y_i`$.
We motivate the commutation relation (12c) from the Yang-Baxter algebra (79) point of view. This also reveals a complete symmetry between the spectral parameters $`z_i`$ and the generators $`y_i`$.
Let us substitute the spectral parameters $`z_i`$ with the affine generators $`y_i`$ in $`L(z_i)`$, and we require that the relation (85) are preserved under the action of the algebra $`L_i`$ on $`\mathrm{\Psi }`$:
$`t_{12}L_1(y_1)L_2(y_2)\mathrm{\Psi }=L_1(y_1)L_2(y_2)\mathrm{\Psi }\overline{t}_{12},`$ (87)
assuming that (85) holds for $`\mathrm{\Psi }`$.
To avoid cumbersome expressions, we use from here the transposed notation $`\overline{a}X`$ for $`X\overline{a}`$. We must then transpose back the final algebraic relations we deduce. In the transposed notations (87) is equivalent to:
$`(t_{12}\overline{t}_{12})L_1(y_1)L_2(y_2)=0,`$ (88)
under the hypothesis that $`t_{12}=\overline{t}_{12}`$ to the right of any expression. Let us for the moment assume that $`\overline{t}_{12}`$ commutes with the symmetrical expressions in $`y_1,y_2`$.
After substituting the expression of $`L_i(y_i)`$ deduced from (84):
$`L_1(y_1)=(yt_{10}y_1t_{10}^1)P_{01},`$ (89)
the term proportional to $`y^0`$ requires that $`\overline{t}_{12}`$ commutes with $`y_1y_2`$, while the term proportional to $`y`$ imposes that:
$`(t_{12}\overline{t}_{12})(y_2t_{01}t_{12}^1+y_1t_{01}^1t_{12})=0,`$ (90)
under the hypothesis that $`t_{01}=\overline{t}_{12}`$ to the right of any expression. This gives:
$`y_2\overline{t}_{12}+(qq^1)y_1\overline{t}_{12}y_1`$ $`=`$ $`0,`$ (91)
$`y_1\overline{t}_{12}y_2\overline{t}_{12}`$ $`=`$ $`0,`$ (92)
which is equivalent to $`y_2\overline{t}_{12}=\overline{t}_{12}^1y_1`$ and implies in particular that $`\overline{t}_{12}`$ commutes with the symmetrical expressions in $`y_1,y_2`$. After transposition, it yields (12c) back.
Alternatively, we can substitute $`z_i`$ for $`y_i`$ in (12c) and verify that the relation is obeyed when we use the expression (86) of $`\overline{t}_i`$.
### D.3 Eigenvalues of the $`y_j`$ in the polynomial case.
We show that the operators $`y_j`$ defined with the polynomial representation 2.2 are triangular matrices. Let us recall the expression of $`y_i`$:
$`y_i`$ $`=`$ $`x_{ii+1}x_{ii+2}\mathrm{}x_{in1}s_ix_{i1}\mathrm{}x_{ii1}`$ (93)
where the operator $`x_{i,j}`$ takes the form for $`i<j`$:
$`x_{ij}=q^1+(qq^1)(1k_{ij}){\displaystyle \frac{z_j}{z_iz_j}},`$ (94)
and the operators $`s_i`$ act as:
$`P(z_1,..,z_i..,z_n)s_i`$ $`=`$ $`cP(z_1..,sz_i,\mathrm{},z_n).`$ (95)
$`x_{12}`$ commutes with $`z_1z_2`$ and with $`z_k`$ for $`k1,2`$. It acts in the following way for on the monomials $`z_1^m,z_2^m`$:
$`z_1^mx_{12}`$ $`=`$ $`q^1z_1^m+(qq^1)(z_1^{m1}z_2+z_1^{m2}z_2^2+\mathrm{}+z_2^m)`$ (96)
$`z_2^mx_{12}`$ $`=`$ $`qz_2^m(qq^1)(z_1^{m1}z_2+z_1^{m2}z_2^2+\mathrm{}+z_1z_2^{m1})`$ (97)
Let us consider which new monomials $`z^\lambda ^{}`$ can appear when one acts with $`x_{12}`$ on the monomial $`z^\lambda `$. First, all the $`\lambda _l^{}`$ for $`l1,2`$ are equal to $`\lambda _l`$. Then, if $`\{\lambda _i^{}\lambda _j^{}\}\{\lambda _i\lambda _j\}`$ with $`\{i,j\}=\{1,2\}`$ and $`\lambda _j^{}\lambda _i^{}`$ , we must have $`\{\lambda _i^{},\lambda _j^{}\}=\{\lambda _ip,\lambda _j+p\}`$ for some integer $`p`$. Finally, if $`\{\lambda _1^{}\lambda _2^{}\}=\{\lambda _1\lambda _2\}`$, the only possibility is that: $`(\lambda _1^{},\lambda _2^{})=(\lambda _2,\lambda _1)`$ with $`\lambda _1>\lambda _2`$.
Let us define an order on the monomials by saying that $`z^\lambda `$ is larger than $`z^\lambda ^{}`$ if either $`\lambda ^{}`$ is obtained from $`\lambda `$ by a sequence of squeezing operations $`\{\lambda _i,\lambda _j\}\{\lambda _i1,\lambda _j+1\}`$ with $`\lambda _i>\lambda _j+1,`$ or $`\lambda ^{}`$ is a permutation of $`\lambda `$ and can be obtained from $`\lambda `$ by a sequence of permutations $`(\lambda _i,\lambda _{i+1})(\lambda _{i+1},\lambda _i)`$ with $`\lambda _i>\lambda _{i+1}`$ . It follows from the above analysis that the action of $`y_j`$ on a monomial produces only monomials which are smaller with respect to this order. Thus the eigenvalues of the operators $`y_j`$ are given by the diagonal elements in the monomial basis.
Given the partition $`\lambda =(\lambda _1,\mathrm{},\lambda _n)`$, the eigenvalues corresponding to the monomials associated to it are all obtained by permutations of the multiplet:
$`(y_j)=c(q)^{1n}(t^{\lambda _j}q^{2(j1)}).`$ (98)
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# 1 Introduction
## 1 Introduction
Why the spacetime is four dimensional is an interesting question in our quest to understand the origin of universe. Nowadays sting theory is a leading candidate of quantum gravity, and IIB matrix model is a candidate of a non-perturbative formulation of string theory . Matrix models may be useful to answer this question if they can consistently describe the dynamics of gauge fields and quantum gravity as expected.
In IIB matrix model, the spacetime may be represented by a particular configuration of matrices. When we are given a matrix configuration as a background field, we can examine its stability by investigating the behavior of the effective action under the change of some parameters of the background. In this framework, the stabilities of fuzzy $`S^2`$, $`S^2\times S^2`$ and $`T^2\times T^2`$ have been studied. This paper extends these investigations into a simple 6 dimensional manifold.
In the previous investigations, the large $`N`$ scaling behavior of the NC gauge theory on these manifolds has been clarified . In supersymmetric models, the effective action scales as $`N^2`$, $`N`$ and $`N^{\frac{4}{3}}`$ in 2, 4 and 6 dimensional manifolds respectively <sup>1</sup><sup>1</sup>1We subtract the universal gauge volume of $`SU(N)/Z_N`$ from the matrix model effective action.. It always scales as $`N^2`$ and the one loop approximation is exact in bosonic models . A major purpose of this paper is to explicitly verify the predicted scaling behavior of the effective action $`(N^{\frac{4}{3}})`$ at the two loop level for a simple 6d fuzzy manifold.
In Gaussian approximations, it has been found that 4 dimensional spacetime tends to minimize the effective action . The advantage of our approach is that the effective actions on homogeneous spaces are much lower than those in the Gaussian approximations which are $`O(N^2)`$. It is because supersymmetry is broken only softly on homogeneous spaces. Although IIB matrix model is supersymmetric, supersymmetry cannot be respected on compact homogenous spacetime. It is precisely why we obtain nonvanishing effective action on these manifolds. It has been a great challenge to explain 4 dimensionality in string theory since the vacuum degenerates as long as it is supersymmetric. In nonperturbative formulation of string theory, the extension of (Euclidean) spacetime may be inevitably finite. Such a view is consistent with finite entropy of de Sitter spacetime in which we are likely to live in. IIB matrix model suggests that the vacuum degeneracy may be lifted due to finite extension of spacetime.
Fuzzy $`S^2`$ can be embedded in Hermitian matrices as
$$A_i=fj_i,$$
(1.1)
where $`j_i`$ are the generators of $`SU(2)`$ with spin $`l`$ and $`f`$ denotes the scale factor of this background. The spin $`l`$ determines the size of $`S^2`$ in our group theoretic construction of $`S^2=SU(2)/U(1)`$ and we can also introduce the scale factor $`f`$ to vary the overall size of $`S^2`$. Thus the parameters of $`S^2`$ in matrix models are a spin $`l`$ and a scale factor $`f`$ for each $`S^2`$.
After setting up our calculation procedure in section 2, we first investigate the stabilities of $`S^2\times S^2\times S^2`$ in three different matrix models with respect to the variation of the spins while assuming the identical scale factor $`f`$ for all $`S^2`$โs. In section 3 and 4 we find that $`S^2\times S^2\times S^2`$ background is metastable and the effective action favors a single large $`S^2`$ in comparison to the remaining $`S^2\times S^2`$ in the matrix models with Myers term. On the other hand, we find that a large $`S^2\times S^2`$ in comparison to the remaining $`S^2`$ is favored in IIB matrix model itself in section 5. These findings are consistent with previous studies . We subsequently investigate $`S^2\times S^2`$ and $`S^2\times S^2\times S^2`$ in IIB matrix model with respect to the variations of the spins and scale factors in section 6 and 7. In this case, we find unstable directions which lower the effective action away from the most symmetric fuzzy $`S^2\times S^2`$ background as suggested by . We conclude in section 8 with discussions.
## 2 IIB type matrix models on fuzzy $`S^2\times S^2\times S^2`$
Since our motivation is to explain the 4 dimensionality of spacetime in superstring theory, it is natural to investigate the stability of fuzzy $`S^2\times S^2\times S^2`$ in IIB matrix model
$$S_{IIB}=\frac{1}{4}Tr[A_\mu ,A_\nu ]^2\frac{1}{2}Tr\overline{\psi }\mathrm{\Gamma }_\mu [A_\mu ,\psi ],$$
(2.1)
where $`A_\mu `$ and $`\psi `$ are $`N\times N`$ Hermitian matrices. We also investigate a deformed model with Myers term
$$S_{Myers}=\frac{i}{3}f_{\mu \nu \rho }Tr[A_\mu ,A_\nu ]A_\rho .$$
(2.2)
It is because fuzzy $`S^2`$, $`S^2\times S^2`$ and $`S^2\times S^2\times S^2`$ become classical solutions after such a deformation since the equations of motion are
$`[A_\mu ,[A_\mu ,A_\nu ]]+\{\overline{\psi },\mathrm{\Gamma }_\nu \psi \}`$ $`=`$ $`\{\begin{array}{cc}if_{\mu \nu \rho }[A_\mu ,A_\rho ]& (\mathrm{with}\mathrm{Myers}\mathrm{term})\\ 0& (\mathrm{without}\mathrm{Myers}\mathrm{term})\end{array},`$ (2.4)
$`\mathrm{\Gamma }_\mu [A_\mu ,\psi ]`$ $`=`$ $`0.`$ (2.5)
This modification does not alter the convergence properties of IIB matrix model while it breaks supersymmetry softly.
After adopting fuzzy $`S^2\times S^2\times S^2`$ as a background field, we separate $`A_\mu `$ and $`\psi `$ into background fields $`p_\mu `$, $`\chi `$ and quantum fluctuations $`a_\mu `$, $`\phi `$.
$`A_\mu `$ $`=`$ $`p_\mu +a_\mu ,`$
$`\psi `$ $`=`$ $`\chi +\phi .`$ (2.6)
Since fuzzy $`S^2`$ can be realized as in (1.1), we take the following $`p_\mu `$ and $`\chi `$ to represent fuzzy $`S^2\times S^2\times S^2`$ as
$`p_\mu `$ $`=`$ $`f\left(\overline{j}_\mu 11\right)1_n(\mu =1,2,3),`$
$`p_\mu `$ $`=`$ $`f\left(1\widehat{j}_\mu 1\right)1_n(\mu =4,5,6),`$
$`p_\mu `$ $`=`$ $`f\left(11\stackrel{~}{j}_\mu \right)1_n(\mu =7,8,9),`$
$`p_0`$ $`=`$ $`0,`$
$`\chi `$ $`=`$ $`0.`$ (2.7)
where $`\overline{j}_\mu ,\widehat{j}_\mu `$ and $`\stackrel{~}{j}_\mu `$ are the angular momentum operators with spin $`l_1`$, $`l_2`$ and $`l_3`$ respectively. $`1_n`$ represents the $`n\times n`$ identity corresponding to $`n`$ coincident fuzzy $`S^2\times S^2\times S^2`$. This background is a classical solution when Myers term is added while it becomes a quantum solution in IIB matrix model at the two loop level when the effective action is stationary with respect to $`f`$.
$`p_\mu `$โs satisfy the $`SU(2)`$ algebras as follows
$`[p_\mu ,p_\nu ]`$ $`=`$ $`if_{\mu \nu \rho }p_\rho ,`$
$`f_{\mu \nu \rho }`$ $`=`$ $`\{\begin{array}{cc}fฯต_{\mu \nu \rho }& (\mu ,\nu ,\rho )(1,2,3)\\ fฯต_{\mu \nu \rho }& (\mu ,\nu ,\rho )(4,5,6)\\ fฯต_{\mu \nu \rho }& (\mu ,\nu ,\rho )(7,8,9)\\ 0& (\mathrm{others})\end{array}.`$ (2.9)
We can construct the Casimir operators in the respective representations as
$`{\displaystyle \underset{\mu =1,2,3}{}}(p_\mu )^2`$ $`=`$ $`f^2l_1(l_1+1),`$
$`{\displaystyle \underset{\mu =4,5,6}{}}(p_\mu )^2`$ $`=`$ $`f^2l_2(l_2+1),`$
$`{\displaystyle \underset{\mu =7,8,9}{}}(p_\mu )^2`$ $`=`$ $`f^2l_3(l_3+1).`$ (2.10)
Since the right hand side of the above equations are the squared radii, the spin and scale factor $`f`$ determine the size of each $`S^2`$. The dimension of the matrices is
$$N=n(2l_1+1)(2l_2+1)(2l_3+1).$$
(2.11)
As we can see, this background represents a simple six dimensional spacetime. The choice of $`S^2`$ facilitates us to carry out detailed analysis of the effective action since the representations of $`SU(2)`$ are well known. By varying the representations (spins) of respective $`S^2`$, we can explore the whole range of the manifolds from the 2 dimensional limit with a single large $`S^2`$ to the 6 dimensional limit with a large $`S^2\times S^2\times S^2`$ with the identical radii. We can thus explore which dimensionality between 2 and 6 is favored by the effective action in this class of backgrounds.
The effective action can be evaluated in a background gauge method by substituting (2) for (2.1) and (2.2). We introduce a gauge fixing term $`S_{gf}`$ and a ghost term $`S_{gh}`$ for gauge fixing. The gauge fixed action is
$`S_{IIB}^{}`$ $``$ $`S_{IIB}+S_{gh}+S_{gf}`$ (2.12)
$`=`$ $`{\displaystyle \frac{1}{4}}Tr[p_\mu ,p_\nu ]^2Tra_\rho ([p_\mu ,[p_\rho ,p_\mu ]])`$
$`+{\displaystyle \frac{1}{2}}Tra_\mu (\delta ^{\mu \nu }P^2+2if_{\mu \nu \rho }P_\rho )a_\nu {\displaystyle \frac{1}{2}}Tr\overline{\phi }\mathrm{\Gamma }^\mu P_\mu \phi +TrbP^2c`$
$`TrP_\mu a_\nu [a_\mu ,a_\nu ]{\displaystyle \frac{1}{4}}Tr[a_\mu ,a_\nu ]^2`$
$`{\displaystyle \frac{1}{2}}Tr\overline{\phi }\mathrm{\Gamma }_\mu [a_\mu ,\phi ]+TrbP_\mu [a_\mu ,c],`$
$`S_{Myers}`$ $`=`$ $`{\displaystyle \frac{i}{3}}f_{\mu \nu \rho }Tr[p_\mu ,p_\nu ]p_\rho +if_{\mu \nu \rho }Tr[p_\mu ,p_\nu ]a_\rho `$ (2.13)
$`if_{\mu \nu \rho }Tra_\mu P_\rho a_\nu +{\displaystyle \frac{i}{3}}f_{\mu \nu \rho }Tr[a_\mu ,a_\nu ]a_\rho .`$
where $`P_\mu X=[p_\mu ,X]`$, and
$`S_{gh}`$ $`=`$ $`TrbP_\mu [p_\mu +a_\mu ,c],`$
$`S_{gf}`$ $`=`$ $`{\displaystyle \frac{1}{2}}Tr(P_\mu a_\mu )^2.`$ (2.14)
## 3 Bosonic matrix model with Myers term
In this section, we investigate a bosonic matrix model with Myers term whose action is
$$S=\frac{1}{4}Tr[A_\mu ,A_\nu ]^2+\frac{i}{3}f_{\mu \nu \rho }Tr[A_\mu ,A_\nu ]A_\rho ,$$
(3.1)
where $`\mu =0,1,\mathrm{},9`$ in correspondence to IIB matrix model. We certainly require 9 matrices to construct fuzzy $`S^2\times S^2\times S^2`$. Although we are eventually interested in IIB matrix model, bosonic models are simple enough to admit exact solutions in the large $`N`$ limit. Our prediction can be verified by comparing with numerical investigations . For this purpose, it is useful to generalize the classical solutions as follows:
$`p_\mu `$ $`=`$ $`\beta \left(\overline{j}_\mu 11\right)1_n(\mu =1,2,3),`$
$`p_\mu `$ $`=`$ $`\beta \left(1\widehat{j}_\mu 1\right)1_n(\mu =4,5,6),`$
$`p_\mu `$ $`=`$ $`\beta \left(11\stackrel{~}{j}_\mu \right)1_n(\mu =7,8,9),`$
$`p_0`$ $`=`$ $`0,`$
$`\chi `$ $`=`$ $`0.`$ (3.2)
The shift of $`\beta `$ away from the classical value of $`f`$ is required to take account of quantum effects. Since the tree and the one loop contributions dominate in bosonic models in the large $`N`$ limit as we shall see, this approach enables us an exact investigation within this class of the backgrounds.
The tree level effective action is
$`\mathrm{\Gamma }_{tree}`$ $`=`$ $`{\displaystyle \frac{1}{4}}Tr[p_\mu ,p_\nu ]^2+{\displaystyle \frac{i}{3}}f_{\mu \nu \rho }Tr[p_\mu ,p_\nu ]p_\rho `$ (3.3)
$``$ $`8N^2\left[{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\beta }{f}}\right)^4{\displaystyle \frac{2}{3}}\left({\displaystyle \frac{\beta }{f}}\right)^3\right]{\displaystyle \frac{f^4}{2^5n^{\frac{2}{3}}N^{\frac{1}{3}}}}(r_1+r_2+r_3).`$
Here the large $`N`$ limit is taken when we proceed from the 1st to 2nd line ($`l_1,l_2,l_31`$) and the following ratios which measure the relative sizes of $`S^2`$โs are introduced
$$r_1=\left(\frac{l_1l_1}{l_2l_3}\right)^{\frac{2}{3}},r_2=\left(\frac{l_2l_2}{l_1l_3}\right)^{\frac{2}{3}},r_3=\left(\frac{l_3l_3}{l_1l_2}\right)^{\frac{2}{3}}.$$
(3.4)
The leading term of the one loop effective action in the large $`N`$ limit can be evaluated as
$`\mathrm{\Gamma }_{1loop}`$ $`=`$ $`{\displaystyle \frac{1}{2}}Tr\mathrm{log}(P^2\delta _{\mu \nu })TrlogP^2`$ (3.5)
$``$ $`8N^2(\mathrm{log}\beta +{\displaystyle \frac{1}{3}}\mathrm{log}{\displaystyle \frac{N}{8n}}`$
$`+{\displaystyle \frac{1}{128}}{\displaystyle _0^4}dXdYdZ\mathrm{log}(r_1X+r_2Y+r_3Z)).`$
The magnitude of the two loop effective action can be estimated by the following power counting argument. The two loop amplitude diverges as the eighth power of the cutoff ($`N^{\frac{8}{3}}`$) while it is suppressed by a power of $`N`$ which is associated with the interaction vertices. Since the one loop amplitude is $`O(N^2)`$, we need to adopt the โt Hoot coupling ($`N^{\frac{1}{3}}/f^4`$) in such a way to make the tree contribution of $`O(N^2)`$ as well. With such a choice, the two loop amplitude scales as $`N^{\frac{4}{3}}`$. Since the tree and the one loop amplitude scale as $`N^2`$, we can safely ignore the two loop amplitude in the large $`N`$ limit. The analogous argument shows that we can ignore all higher loop contributions and the one loop effective action becomes exact in the large $`N`$ limit.
In this way, we obtain the effective action in the large $`N`$ limit as
$`\mathrm{\Gamma }`$ $`=`$ $`8N^2([{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\beta }{f}}\right)^4{\displaystyle \frac{2}{3}}\left({\displaystyle \frac{\beta }{f}}\right)^3]{\displaystyle \frac{1}{2^5\lambda _1^2}}(r_1+r_2+r_3)`$ (3.6)
$`+\mathrm{log}\beta +{\displaystyle \frac{1}{3}}\mathrm{log}{\displaystyle \frac{N}{8n}}+{\displaystyle \frac{1}{128}}{\displaystyle _0^4}dXdYdZ\mathrm{log}(r_1X+r_2Y+r_3Z)),`$
where $`\lambda _1`$ is the โt Hooft coupling:
$$\lambda _1^2=\frac{n^{\frac{2}{3}}N^{\frac{1}{3}}}{f^4}.$$
(3.7)
To minimize the effective action, we have to solve $`\frac{\mathrm{\Gamma }}{\beta }=0`$. This condition determines the scale factor $`\beta `$ as
$`{\displaystyle \frac{\beta }{f}}`$ $`=`$ $`{\displaystyle \frac{1}{4}}+{\displaystyle \frac{1}{2}}\sqrt{{\displaystyle \frac{1}{4}}+g(A)}+{\displaystyle \frac{1}{2}}\sqrt{{\displaystyle \frac{1}{2}}g(A)+{\displaystyle \frac{1}{4\sqrt{\frac{1}{4}+g(A)}}}},`$ (3.8)
where
$`g(A)`$ $`=`$ $`{\displaystyle \frac{2^{\frac{5}{3}}A}{3^{\frac{1}{3}}\left(9A+\sqrt{3}\sqrt{27A^2128A^3}\right)^{\frac{1}{3}}}}+{\displaystyle \frac{\left(9A+\sqrt{3}\sqrt{27A^2128A^3}\right)^{\frac{1}{3}}}{6^{\frac{2}{3}}}},`$ (3.9)
and $`A=2^5\lambda _1^2`$. We find that this solution exists in the weak coupling regime where
$$\lambda _1^2=\frac{n^{\frac{2}{3}}N^{\frac{1}{3}}}{f^4}<\frac{3^3}{2^{12}}(r_1+r_2+r_3)0.0198\mathrm{a}tr_1=r_2=r_3=1.$$
(3.10)
Beyond this critical point, the fuzzy $`S^2\times S^2\times S^2`$ background no longer exists. Just like , this point separates the background distributions between the collapsed phase and the fuzzy $`S^2\times S^2\times S^2`$ phase.
After plugging (3.6) into (3.8), we can estimate the effective action by performing the triple integrals numerically for each $`\lambda _1`$. Fig.1 shows $`\mathrm{\Gamma }/8N^2`$ against $`t=l_1/l_3=l_2/l_3`$ for $`N=8\times 10^6`$ and $`\lambda _1^2=0.0197,0.005,0.001`$. $`t=1`$ corresponds to the most symmetric fuzzy $`S^2\times S^2\times S^2`$ background. As $`t\mathrm{}`$, the background approaches fuzzy $`S^2\times S^2`$ while it approaches fuzzy $`S^2`$ as $`t0`$.
From Fig.1 we can see that fuzzy $`S^2\times S^2\times S^2`$ background is not stable. To explore whether it tends to decay into $`S^2\times S^2`$ or $`S^2`$, let us examine the behavior of the effective action in more details. When we fix $`l_1/l_3=100`$ and vary $`l_2/l_3`$ between 1 to 100, the action behaves as Fig.2. Here, $`l_2/l_3=100`$ and $`l_2/l_3=1`$ in Fig.2 correspond to $`t=100`$ and $`t=0.01`$ in Fig.1 respectively because of the equivalence of the three $`S^2`$โs. As we observe $`\mathrm{\Gamma }`$ is smooth with respect to $`l_1/l_3`$ and $`l_2/l_3`$, and does not develop a local minimum, we can convince ourselves that $`S^2`$ is favored in this example. The situation like this holds for the other combinations of $`l_1/l_3`$ and $`l_2/l_3`$ when one of them is large enough. Therefore we conclude that fuzzy $`S^2\times S^2\times S^2`$ background is not stable in the action (3.1), and it decays toward fuzzy $`S^2`$.
## 4 Deformed IIB matrix model with Myers term
In this section, we study a deformed IIB matrix model with Myers term whose action is
$$S=\frac{1}{4}Tr[A_\mu ,A_\nu ]^2\frac{1}{2}Tr\overline{\psi }\mathrm{\Gamma }_\mu [A_\mu ,\psi ]+\frac{i}{3}f_{\mu \nu \rho }Tr[A_\mu ,A_\nu ]A_\rho .$$
(4.1)
We investigate the stability of the fuzzy $`S^2\times S^2\times S^2`$ background (2) which is a classical solution in this action <sup>2</sup><sup>2</sup>2We need not generalize $`f`$ into $`\beta `$ here since there are no tadpoles at the one loop level ..
The effective action can be evaluated by extending the procedures in . The results in the large N limit are
$`\mathrm{\Gamma }_{tree}`$ $``$ $`n^{\frac{2}{3}}N^{\frac{4}{3}}{\displaystyle \frac{1}{24}}{\displaystyle \frac{f^4N^{\frac{1}{3}}}{n^{\frac{4}{3}}}}(r_1+r_2+r_3),`$
$`\mathrm{\Gamma }_{1loop}`$ $``$ $`n^{\frac{2}{3}}N^{\frac{4}{3}}{\displaystyle \frac{1}{8}}{\displaystyle _0^4}๐X{\displaystyle _0^4}๐L{\displaystyle _0^4}๐A{\displaystyle \frac{1}{r_1X+r_2L+r_3A}},`$
$`\mathrm{\Gamma }_{2loop}`$ $``$ $`n^{\frac{2}{3}}N^{\frac{4}{3}}{\displaystyle \frac{n^{\frac{4}{3}}}{f^4N^{\frac{1}{3}}}}\left({\displaystyle \frac{3}{2}}f_3+4f_4\right),`$ (4.2)
$`f_3`$ and $`f_4`$ are defined as the following multiple integrals
$`f_3`$ $`=`$ $`{\displaystyle _0^4}{\displaystyle \frac{dXdYdZdLdMdNdAdBdC}{(r_1X+r_2L+r_3A)(r_1Y+r_2M+r_3B)(r_1Z+r_2N+r_3C)}}`$
$`\times W(X,Y,Z)W(L,M,N)W(A,B,C),`$
$`f_4`$ $`=`$ $`{\displaystyle _0^4}{\displaystyle \frac{dXdYdZdLdMdNdAdBdC}{(r_1X+r_2L+r_3A)^2(r_1Y+r_2M+r_3B)^2(r_1Z+r_2N+r_3C)}}`$ (4.3)
$`\times [r_1^2XY+r_2^2LM+r_3^2AB+{\displaystyle \frac{1}{2r_3}}(ZXY)(NLM)`$
$`\times +{\displaystyle \frac{1}{2r_2}}(ZXY)(CAB)+{\displaystyle \frac{1}{2r_1}}(NLM)(CAB)]`$
$`\times W(X,Y,Z)W(L,M,N)W(A,B,C),`$
where $`W(X,Y,Z)`$ is the asymptotic formula of Wignerโs $`6j`$ symbols which appear in the interaction vertices:
$`l_1^3\left\{\begin{array}{ccc}j_1& j_2& j_3\\ l_1& l_1& l_1\end{array}\right\}^2`$ $``$ $`W(j_1^2,j_2^2,j_3^2)`$ (4.6)
$`=`$ $`{\displaystyle \frac{1}{2\pi \sqrt{\frac{YZ(4Y)(4Z)}{4}\left(X\frac{2Y+2ZYZ}{2}\right)^2}}}.`$ (4.7)
This approximation is valid in the uniformly large angular momentum regime. Since the effective action is highly divergent, we argue that this approximation is exact in the large $`N`$ limit.
The effective action at the 2-loop level is
$$\mathrm{\Gamma }=n^{\frac{2}{3}}N^{\frac{4}{3}}\left[\frac{r_1+r_2+r_3}{24\lambda _2^2}+\frac{1}{8}_0^4\frac{dXdLdA}{r_1X+r_2L+r_3A}\lambda _2^2\left(\frac{3}{2}f_3+4f_4\right)\right].$$
(4.8)
It is indeed of $`O(N^{\frac{4}{3}})`$ after choosing the โt Hooft coupling $`\lambda _2`$:
$$\lambda _2^2=\frac{n^{\frac{4}{3}}}{f^4N^{\frac{1}{3}}}.$$
(4.9)
In the same way as in the previous section, we have estimated this effective action numerically. From Fig.3 we can observe again that the fuzzy $`S^2\times S^2\times S^2`$ background is not stable. The same analysis in the previous section can be used here to determine into which configuration it tends to decay. From such an analysis we conclude that the fuzzy $`S^2\times S^2\times S^2`$ background decays toward fuzzy $`S^2`$. It is interesting to observe that $`S^2\times S^2\times S^2`$ becomes a local minimum of the effective action at $`t=1`$ when the coupling $`\lambda _2`$ is strong enough.
## 5 IIB matrix model analysis with respect to spins
After investigating bosonic and supersymmetric models with Myers term, we evaluate the effective action for IIB matrix model (2.1) with the background (2). The results are <sup>3</sup><sup>3</sup>3The two loop amplitude $`\mathrm{\Gamma }_{2loop}`$ is evaluated in the Appendix with generic scale factors.
$`\mathrm{\Gamma }_{tree}`$ $``$ $`n^{\frac{2}{3}}N^{\frac{4}{3}}{\displaystyle \frac{1}{8}}{\displaystyle \frac{f^4N^{\frac{1}{3}}}{n^{\frac{4}{3}}}}(r_1+r_2+r_3),`$
$`\mathrm{\Gamma }_{1loop}`$ $``$ $`O\left(N^{\frac{2}{3}}\right),`$
$`\mathrm{\Gamma }_{2loop}`$ $``$ $`n^{\frac{2}{3}}N^{\frac{4}{3}}{\displaystyle \frac{n^{\frac{4}{3}}}{f^4N^{\frac{1}{3}}}}{\displaystyle \frac{1}{2}}f_3.`$ (5.1)
We can explicitly check that the effective action scales as $`N^{\frac{4}{3}}`$ at the two loop level with the following choice of the โt Hoot coupling
$`\mathrm{\Gamma }`$ $`=`$ $`n^{\frac{2}{3}}N^{\frac{4}{3}}\left({\displaystyle \frac{1}{8\lambda _2^2}}(r_1+r_2+r_3)+{\displaystyle \frac{\lambda _2^2}{2}}f_3\right),`$
$`\lambda _2^2`$ $`=`$ $`{\displaystyle \frac{n^{\frac{4}{3}}}{f^4N^{\frac{1}{3}}}}.`$ (5.2)
Since the โt Hoot coupling which is set by the overall scale of the background becomes a dynamical variable, we can minimize the effective action with respect to it for fixed representations:
$$\mathrm{\Gamma }_{min}=2\sqrt{\mathrm{\Gamma }_{tree}\mathrm{\Gamma }_{2loop}}.$$
(5.3)
We can now explore the minimum of the effective action (5.3) with respect to the representations $`(l_1,l_2,l_3)`$. Fig.4 shows $`\mathrm{\Gamma }_{min}/n^{\frac{3}{2}}N^{\frac{4}{3}}`$ against $`t=l_1/l_3=l_2/l_3`$. We note that the action is decreasing faster in the $`t>1`$ region than $`t<1`$ region. We may thus conclude that fuzzy $`S^2\times S^2\times S^2`$ is not stable and it tends to decay toward fuzzy $`S^2\times S^2`$. This result is consistent with the previous investigations .
We can further demonstrate that the action (5) reduces to that of $`S^2\times S^2`$ when we take a limit $`l_1,l_2l_3`$. It is because we can reexpress (5) as
$`\mathrm{\Gamma }_{tree}^{\prime \prime }`$ $`=`$ $`{\displaystyle \frac{N}{8\lambda _3^2}}\left(r+1/r+R\right),`$
$`\mathrm{\Gamma }_{2loop}^{\prime \prime }`$ $`=`$ $`2N\lambda _3^2l_3^2{\displaystyle _0^4}๐X๐Y๐Z๐L๐M๐N๐A๐B๐C`$
$`\times {\displaystyle \frac{W(X,Y,Z)W(L,M,N)W(A,B,C)}{(rX+r^1L+RA)(rY+r^1M+RB)(rZ+r^1N+RC)}},`$
$`N`$ $``$ $`2l_12l_22l_3,\lambda _3^2={\displaystyle \frac{1}{f^42l_12l_2}},R={\displaystyle \frac{l_3^2}{l_1l_2}},r={\displaystyle \frac{l_1}{l_2}}.`$ (5.4)
Here, we set $`n=1`$ for simplicity.
Fig.5 shows $`\mathrm{\Gamma }_{min}^{\prime \prime }/N=2\sqrt{\mathrm{\Gamma }_{tree}^{\prime \prime }\mathrm{\Gamma }_{2loop}^{\prime \prime }}/N`$ against $`t=l_1/l_3=l_2/l_3`$ for $`N=8\times 10^9`$. It demonstrates that the effective action scales in a 4d fashion as the background approaches $`S^2\times S^2`$.
## 6 Stability of fuzzy $`S^2\times S^2`$ background
In this section, we investigate the stability of fuzzy $`S^2\times S^2`$ background in detail with respect to the scale factors in addition to the change of the representations (spins). Our investigation is motivated by an analogous study for the fuzzy torus case .
We consider the background of the following type:
$`p_\mu `$ $`=`$ $`f_1\left(\overline{j}_\mu 1\right)(\mu =1,2,3),`$
$`p_\mu `$ $`=`$ $`f_2\left(1\widehat{j}_\mu \right)(\mu =4,5,6),`$
$`p_\mu `$ $`=`$ $`0(\mu =7,8,9,0),`$
$`\chi `$ $`=`$ $`0.`$ (6.1)
Here we set $`n=1`$ for simplicity. It generalizes the background in where the identical scale factor $`f_1=f_2`$ is assumed.
The effective actions are evaluated as <sup>4</sup><sup>4</sup>4It can again be read off from $`\mathrm{\Gamma }_{2loop}`$ in the Appendix.
$`\mathrm{\Gamma }_{}^{}{}_{tree}{}^{}`$ $`=`$ $`{\displaystyle \frac{N^{}}{8\lambda _{}^{}{}_{}{}^{2}}}\left(rq^2+{\displaystyle \frac{1}{rq^2}}\right),`$
$`\mathrm{\Gamma }_{}^{}{}_{1loop}{}^{}`$ $`=`$ $`O(\mathrm{log}N^{}),`$
$`\mathrm{\Gamma }_{}^{}{}_{2loop}{}^{}`$ $`=`$ $`4N^{}\lambda _{}^{}{}_{}{}^{2}{\displaystyle _0^4}๐X๐Y๐Z๐L๐M๐N`$ (6.2)
$`\times {\displaystyle \frac{(rq^2X+\frac{1}{rq^2}L)}{(rqX+\frac{1}{rq}L)^2(rqY+\frac{1}{rq}M)(rqZ+\frac{1}{rq}N)}}`$
$`\times W(X,Y,Z)W(L,M,N),`$
where
$$N^{}=2l_12l_2,\lambda _{}^{}{}_{}{}^{2}=\frac{1}{f_1^2f_2^2N^{}},q=\frac{f_1}{f_2},r=\frac{l_1}{l_2}.$$
(6.3)
Fig.6 and Fig.7 show $`\mathrm{\Gamma }_{min}^{}/N^{}=2\sqrt{\mathrm{\Gamma }_{tree}^{}\mathrm{\Gamma }_{2loop}^{}}/N^{}`$ against $`r`$ and $`q`$. The points represented by squares possess smaller effective actions than that of the most symmetric point $`r=q=1`$. Furthermore $`\sqrt{rq}`$ (the scale ratios of the two $`S^2`$) also decreases in this domain. Therefore we find that fuzzy $`S^2\times S^2`$ background is not stable when the both spins and scale factors are allowed to change at the two loop level.
## 7 Local stability of the backgrounds
In the previous sections, we have investigated the stability of the backgrounds globally by performing the multiple integrals numerically. In this section we also investigate the local stability of the most symmetric background with respect to the small variations of the spins and scale factors.
In $`S^2\times S^2`$ case, the minimum of the effective action (6) under the variations of $`r=1+\delta ,q=1+ฯต,(\delta ,ฯต1)`$ becomes
$`{\displaystyle \frac{\mathrm{\Gamma }_{min}^{}}{N^{}}}`$ $`=`$ $`{\displaystyle \frac{2\sqrt{\mathrm{\Gamma }_{tree}^{}\mathrm{\Gamma }_{2loop}^{}}}{N^{}}}{\displaystyle \frac{G}{4}},`$
$`G`$ $`=`$ $`{\displaystyle \frac{(X+q^4r^2L)(1+q^4r^2)r^2}{(X+q^2r^2L)^2(Y+q^2r^2M)(Z+q^2r^2N)}}`$ (7.1)
$``$ $`2F^{}+({\displaystyle \frac{7}{2}}F^{}+6|C^{}|^2+6|\epsilon |^2)\delta ^2`$
$`+2(|C^{}|^2+|\epsilon ^{}|^2)ฯต^2+(4F^{}+8|C^{}|^2+8|\epsilon ^{}|^2)\delta ฯต,`$
where
$`F^{}`$ $``$ $`{\displaystyle \frac{1}{(X+L)(Y+M)(Z+N)}},`$
$`|C^{}|^2`$ $``$ $`{\displaystyle \frac{L^2}{(X+L)^3(Y+M)(Z+N)}},`$
$`|\epsilon ^{}|^2`$ $``$ $`{\displaystyle \frac{LM}{(X+L)^2(Y+M)^2(Z+N)}},`$
$`\mathrm{}`$ $``$ $`{\displaystyle _0^4}๐X๐Y๐Z๐L๐M๐N\mathrm{}W(X,Y,Z)W(L,M,N).`$ (7.2)
We can estimate $`F^{},|C^{}|^2,|\epsilon ^{}|^2`$ numerically as follows
$`F^{}`$ $`=`$ $`3.263930\pm 0.19\times 10^4,`$
$`|C^{}|^2`$ $`=`$ $`1.056778\pm 0.94\times 10^5,`$
$`|\epsilon ^{}|^2`$ $`=`$ $`0.8630541\pm 0.68\times 10^5.`$ (7.3)
Let us consider the following ratio:
$$\frac{G}{F^{}}=2+0.0292\delta ^2+1.18ฯต^2+0.706ฯต\delta .$$
(7.4)
To explore the stability of the background, we evaluate the eigenvalues of the quadratic forms in $`\delta `$ and $`ฯต`$ above. The eigenvalues are
$$1.28,0.07.$$
(7.5)
Since we find a negative eigenvalue, we conclude that $`S^2\times S^2`$ is not stable in such a direction around the most symmetric point.
The same analysis can be applied to $`S^2\times S^2\times S^2`$. In this case, a relevant background corresponding to (6) is
$`p_\mu `$ $`=`$ $`f_1\left(\overline{j}_\mu 11\right)1_n(\mu =1,2,3),`$
$`p_\mu `$ $`=`$ $`f_2\left(1\widehat{j}_\mu 1\right)1_n(\mu =4,5,6),`$
$`p_\mu `$ $`=`$ $`f_3\left(11\stackrel{~}{j}_\mu \right)1_n(\mu =7,8,9),`$
$`p_0`$ $`=`$ $`0,`$
$`\chi `$ $`=`$ $`0.`$ (7.6)
We evaluate the two loop effective action for this generic background in the Appendix. The minimum of the effective action with respect to the โt Hooft coupling is
$`\mathrm{\Gamma }_{min}`$ $`=`$ $`{\displaystyle \frac{N^{4/3}H^{1/2}}{2^{1/2}}},`$
$`H`$ $`=`$ $`{\displaystyle \frac{(X+q_2^4r_2^2L+q_3^4r_3^2A)(1+q_2^4r_2^2+q_3^4r_3^2)(r_2^2r_3^2)^{\frac{2}{3}}}{(X+q_2^4r_2^2L+q_3^4r_3^2A)^2(Y+q_2^2r_2^2M+q_3^2r_3^2B)(Z+q_2^2r_2^2N+q_3^2r_3^2C)}},`$
$`\mathrm{}`$ $``$ $`{\displaystyle _0^4}๐X๐Y๐Z๐L๐M๐N๐A๐B๐C\mathrm{}`$ (7.7)
$`\times W(X,Y,Z)W(L,M,N)W(A,B,C),`$
$`r_2={\displaystyle \frac{l_2}{l_1}},r_3={\displaystyle \frac{l_3}{l_1}},q_2={\displaystyle \frac{f_2}{f_1}},q_3={\displaystyle \frac{f_3}{f_1}}.`$
The functional dependence of $`H`$ on $`r_i=1+\delta _i,q_i=1+ฯต_i`$ is
$`H`$ $`=`$ $`3F+(9|C|^2+9|\epsilon |^2{\displaystyle \frac{8}{3}}F)(\delta _2^2+\delta _3^2)({\displaystyle \frac{17}{3}}F9|C|^2)\delta _2\delta _3`$ (7.8)
$`+(9|C|^2+9|\epsilon |^2{\displaystyle \frac{3}{2}}F)(ฯต_2^2+ฯต_3^2)+(F3|C|^2)ฯต_2ฯต_3`$
$`+(12|C|^2+12|\epsilon |^2{\displaystyle \frac{8}{3}}F)(\delta _2ฯต_2+\delta _3ฯต_3)`$
$`+({\displaystyle \frac{3}{2}}F6|C|^2+12|\epsilon |^2)(\delta _2ฯต_3+\delta _3ฯต_2),`$
where
$`F`$ $``$ $`{\displaystyle \frac{1}{(X+L+A)(Y+M+B)(Z+N+C)}},`$
$`|C|^2`$ $``$ $`{\displaystyle \frac{A^2}{(X+L+A)^3(Y+M+B)(Z+N+C)}},`$
$`|\epsilon |^2`$ $``$ $`{\displaystyle \frac{AB}{(X+L+A)^2(Y+M+B)^2(Z+N+C)}}.`$ (7.9)
We can estimated $`F,|C|^2,|\epsilon |^2`$ numerically as follows
$`F`$ $`=`$ $`3.993\pm 0.25\times 10^2,`$
$`|C|^2`$ $`=`$ $`0.6032\pm 0.46\times 10^3,`$
$`|\epsilon |^2`$ $`=`$ $`0.4656\pm 0.11\times 10^2.`$ (7.10)
It is again useful to consider the following ratio:
$`{\displaystyle \frac{H}{F}}`$ $`=`$ $`30.26(\delta _2^2+\delta _3^2)4.3\delta _2\delta _3`$ (7.11)
$`+0.91(ฯต_2^2+ฯต_3^2)0.55ฯต_2ฯต_3`$
$`+0.55(\delta _2ฯต_2+\delta _3ฯต_3)1.0(\delta _2ฯต_3+\delta _3ฯต_2).`$
The eigenvalues of this quadratic form are
$$2.43,2.39,0.69,0.65.$$
(7.12)
The existence of the negative eigenvalue implies the instability of this background. If we consider the $`\delta _2,\delta _3`$ subspace, the eigenvalues are
$$2.41,1.89$$
(7.13)
indicating the instability of this background in agreement with section 5.
## 8 Conclusions
In this paper we have investigated the effective action of matrix models with $`S^2\times S^2\times S^2`$ backgrounds at the two loop level. This class of 6 dimensional manifolds can be constructed by using $`SU(2)`$ algebra which facilitates us to evaluate the effective action. We can change the size of each $`S^2`$ by choosing different representations (spins) of $`SU(2)`$. Therefore we can probe manifolds of different dimensionality such as $`S^2`$ (2d), $`S^2\times S^2`$ (4d) and $`S^2\times S^2\times S^2`$ (6d) by collapsing some of $`S^2`$โs. Since the background with the smallest effective action is most likely to be realized in a particular model, this investigation may shed light why our spacetime is 4 dimensional in IIB matrix model context.
In the previous investigations, the large $`N`$ scaling behavior of the NC gauge theory on these manifolds has been clarified. In supersymmetric models, the effective action scales as $`N^2`$, $`N`$ and $`N^{\frac{4}{3}}`$ in 2, 4 and 6 dimensional manifolds respectively. It is always $`O(N^2)`$ and the one loop approximation is exact in bosonic models. We have indeed verified this scaling behavior for 6 dimensional spacetime at the two loop level. With the presence of Myers term, $`S^2`$ configuration is favored since the effective action can become negative. On the other hand, 4 dimensional spacetime is favored in IIB matrix model since the effective action is positive definite.
In accord with these expectations we find that the fuzzy $`S^2\times S^2\times S^2`$ background is not stable when we vary the ratios of the spins in the following matrix models
* 10d bosonic matrix model with Myers term,
* IIB matrix model with Myers term,
* IIB matrix model.
In the first two matrix models with Myers term, $`S^2\times S^2\times S^2`$ tends to decay toward $`S^2`$. On the other hand $`S^2\times S^2\times S^2`$ tends to decay $`S^2\times S^2`$ in the last one. These results are consistent with the previous works .
We have further investigated the effective action around the symmetric $`S^2\times S^2`$ in detail. Under the change of the ratios of the scale factors in addition to the spins, we find there are unstable directions of the $`S^2\times S^2`$ background in IIB matrix model.
This instability does not imply that it eventually decays into $`S^2`$ since the effective action is $`O(N^2)`$ in such a limit. However this argument cannot be verified at the two loop level since the effective action for $`S^2`$ behaves as follows up to the two loop level
$$aN^3f^4+\frac{b}{f^4N},$$
(8.14)
where $`a,b`$ are $`O(1)`$. It is therefore possible to make it $`O(N)`$ by choosing $`1/f^4N^2`$ to be $`O(1)`$ which is the same order with 4d manifolds. However we argue that this is a two loop artifact since the $`n`$ loop contribution can be estimated as $`(1/f^4N)^{n1}`$ and hence we need to adopt $`1/f^4N`$ as the โt Hooft coupling. Nevertheless we cannot rule out the possible existence of instability of a simple product space of $`S^2\times S^2`$. In this respect a more symmetric 4d spacetime such as $`CP^2=SU(3)/U(2)`$ is interesting and it may not suffer from the instability found here since $`SU(3)`$ symmetry permits only the over all scale factor of the background . It is also likely that a manifold with higher symmetry may lower the effective action of IIB matrix model.
Acknowledgments
This work is supported in part by the Grant-in-Aid for Scientific Research from the Ministry of Education, Science and Culture of Japan. We thank H. Aoki, T. Azuma, S. Bal , S. Iso, H. Kawai, K. Nagao, J. Nishimura and Y. Takayama for discussions.
## Appendix A
In this appendix, we evaluate the two loop effective action of IIB matrix model in the $`S^2\times S^2\times S^2`$ background (7). Our calculation which is an extension of those in incorporates the asymmetric scale factors. The result in the $`S^2\times S^2`$ background (6) can be obtained by shrinking one of $`S^2`$โs.
We expand quantum fluctuations in terms of the tensor product of the matrix spherical harmonics:
$`a^\mu `$ $`=`$ $`{\displaystyle \underset{jmpqst}{}}a_{jmpqst}^\mu (Y_{jm}Y_{pq}Y_{st}),`$
$`\phi `$ $`=`$ $`{\displaystyle \underset{jmpqst}{}}\phi _{jmpqst}(Y_{jm}Y_{pq}Y_{st}),`$
$`b`$ $`=`$ $`{\displaystyle \underset{jmpqst}{}}b_{jmpqst}(Y_{jm}Y_{pq}Y_{st}),`$
$`c`$ $`=`$ $`{\displaystyle \underset{jmpqst}{}}c_{jmpqst}(Y_{jm}Y_{pq}Y_{st}).`$ (A.1)
Here the sums of $`j`$, $`p`$ and $`s`$ run up to $`2l_1`$, $`2l_2`$ and $`2l_3`$ respectively. Then the propagators are derived from the kinetic terms of (2.12):
$`a^\mu a^\nu `$ $`=`$ $`{\displaystyle \underset{jmpqst}{}}\left(P^2\delta _{\mu \nu }+2if_{\mu \nu \rho }P^\rho \right)^1(Y_{jm}Y_{pq}Y_{st})(Y_{jm}^{}Y_{pq}^{}Y_{st}^{}),`$
$`\phi \overline{\phi }`$ $`=`$ $`{\displaystyle \underset{jmpqst}{}}\left(\mathrm{\Gamma }_\mu P_\mu \right)^1(Y_{jm}Y_{pq}Y_{st})(Y_{jm}^{}Y_{pq}^{}Y_{st}^{}),`$
$`<cb>`$ $`=`$ $`{\displaystyle \underset{jmpqst}{}}{\displaystyle \frac{1}{P^2}}(Y_{jm}Y_{pq}Y_{st})(Y_{jm}^{}Y_{pq}^{}Y_{st}^{}).`$ (A.2)
We exclude the singlet state $`j=p=s=0`$ in the propagators. To calculate the leading contributions in the large $`N`$ limit, we expand the boson and the fermion propagators as
$`\left(P^2\delta _{\mu \nu }+2if_{\mu \nu \rho }P^\rho \right)^1`$ $``$ $`{\displaystyle \frac{\delta _{\mu \nu }}{P^2}}2i{\displaystyle \frac{f_{\mu \nu \rho }P^\rho }{P^4}}+4{\displaystyle \frac{I_{\mu \nu }(P)}{P^6}},`$
$`\left(\mathrm{\Gamma }_\mu P_\mu \right)^1`$ $``$ $`{\displaystyle \frac{\mathrm{\Gamma }^\mu P_\mu }{P^2}}+{\displaystyle \frac{i}{2}}{\displaystyle \frac{f_{\mu \nu \sigma }\mathrm{\Gamma }^{\mu \nu \rho }P_\sigma P_\rho }{P^4}}{\displaystyle \frac{\mathrm{\Gamma }f^2P}{P^4}}`$ (A.3)
$`+{\displaystyle \frac{Pf^2PP^\mu \mathrm{\Gamma }_\mu }{P^6}}.`$
We have introduced the following tensor
$$I_{\mu \nu }(\overline{\delta }_{\mu \nu }\overline{P}^2\overline{P}_{\mu \nu })f_1^2+(\widehat{\delta }_{\mu \nu }\widehat{P}^2\widehat{P}_{\mu \nu })f_2^2+(\stackrel{~}{\delta }_{\mu \nu }\stackrel{~}{P}^2\stackrel{~}{P}_{\mu \nu })f_3^2.$$
(A.4)
The symbols $`\overline{}`$, $`\widehat{}`$ and $`\stackrel{~}{}`$ denote the sub-spaces whose Lorentz indices $`\mu `$ run over $`(1,2,3)`$, $`(4,5,6)`$ and $`(7,8,9)`$ respectively. We also introduce
$$Pf^2Pf_1^2\overline{P}^2+f_2^2\widehat{P}^2+f_3^2\stackrel{~}{P}^2.$$
(A.5)
Using these propagators, we can calculate the contributions to the two loop effective action from various interaction vertices as follows.
4-gauge boson vertex is
$$V_4=\frac{1}{4}Tr[a_\mu ,a_\nu ]^2.$$
(A.6)
The leading contribution to the two loop effective action is
$$<V_4>=\frac{1}{P_1^2P_2^2}\left\{45+6\frac{P_3f^2P_3}{P_1^2P_2^2}12\frac{P_2f^2P_2}{P_1^2P_2^2}72\frac{P_1f^2P_1}{P_1^4}\right\}_P.$$
(A.7)
We introduce the wave functions and averages as
$`\mathrm{\Psi }_{123}`$ $``$ $`Tr(Y_{j_1m_1}Y_{j_2m_2}Y_{j_3m_3})Tr(Y_{p_1q_1}Y_{p_2q_2}Y_{p_3q_3})Tr(Y_{s_1t_1}Y_{s_2t_2}Y_{s_3t_3}),`$
$`X_P`$ $``$ $`{\displaystyle \underset{j_i,p_i,s_i,m_i,q_i,t_i}{}}\mathrm{\Psi }_{123}^{}X\mathrm{\Psi }_{123},`$
$`P_i^\mu (Y_{j_im_i}Y_{p_iq_i}Y_{s_it_i})`$ $``$ $`[p_\mu ,Y_{j_im_i}Y_{p_iq_i}Y_{s_it_i}].`$ (A.8)
We define following functions:
$`F_1`$ $`=`$ $`{\displaystyle \frac{1}{P_1^4P_2^4}}_P,`$
$`\stackrel{~}{g}_1`$ $`=`$ $`{\displaystyle \frac{P_2f^2P_2}{P_1^4P_2^4}}_P,`$
$`g_1`$ $`=`$ $`{\displaystyle \frac{P_1f^2P_1}{P_1^6P_2^2}}_P,`$
$`g_2`$ $`=`$ $`{\displaystyle \frac{P_3f^2P_3}{P_1^4P_2^4}}_P.`$ (A.9)
Then
$$<V_4>=45F_112\stackrel{~}{g}_172g_1+6g_2.$$
(A.10)
Ghost vertex is
$$V_g=Trb[p_\mu ,[a_\mu ,c]].$$
(A.11)
Their contribution is
$$\frac{1}{2}<V_gV_g>=F_2+4H_2.$$
(A.12)
Here
$`F_2`$ $`=`$ $`{\displaystyle \frac{P_2P_3}{P_1^2P_2^2P_3^2}}_P,`$
$`H_2`$ $`=`$ $`{\displaystyle \frac{P_2I(1)P_3}{P_1^6P_2^2P_3^2}}_P,`$ (A.13)
and
$$P_iI(j)P_kP_i^\mu I_{\mu \nu }(P_j)P_k^\nu .$$
(A.14)
3-gauge boson vertex is
$$V_3=TrP_\mu a_\nu [a_\mu ,a_\nu ].$$
(A.15)
Their contribution is
$`{\displaystyle \frac{1}{2}}<V_3V_3>`$ $`=`$ $`9F_19F_2+12F_3+8g_1^{}4\stackrel{~}{g}_1^{}+2g_2`$ (A.16)
$`+32H_136H_216H_3+12H_44H_5.`$
Newly introduced functions are defined as
$`F_3`$ $`=`$ $`{\displaystyle \frac{P_1f^2P_1}{P_1^4P_2^2P_3^2}}_P,`$
$`g_1^{}`$ $`=`$ $`g_1{\displaystyle \frac{1}{N}}{\displaystyle \underset{j,p,s}{}}(2j+1)(2p+1)(2s+1)`$
$`\times {\displaystyle \frac{f_1^4j(j+1)+f_2^4p(p+1)+f_3^4s(s+1)}{\left[f_1^2j(j+1)+f_2^2p(p+1)+f_3^2s(s+1)\right]^4}},`$
$`\stackrel{~}{g}_1^{}`$ $`=`$ $`\stackrel{~}{g}_1{\displaystyle \frac{1}{N}}{\displaystyle \underset{j,p,s}{}}(2j+1)(2p+1)(2s+1)`$
$`\times {\displaystyle \frac{f_1^4j(j+1)+f_2^4p(p+1)+f_3^4s(s+1)}{\left[f_1^2j(j+1)+f_2^2p(p+1)+f_3^2s(s+1)\right]^4}},`$
$`H_1`$ $`=`$ $`{\displaystyle \frac{P_1I(2)P_1}{P_1^2P_2^6P_3^2}}_P,`$
$`H_3`$ $`=`$ $`{\displaystyle \frac{P_2I(1)P_3}{P_1^4P_2^4P_3^2}}_P,`$
$`H_4`$ $`=`$ $`{\displaystyle \frac{P_1I(2)P_1}{P_1^4P_2^4P_3^2}}_P,`$
$`H_5`$ $`=`$ $`{\displaystyle \frac{P_2I(1)P_3}{P_1^2P_2^4P_3^4}}_P.`$ (A.17)
Fermion vertex is
$$V_f=\frac{1}{2}Tr\overline{\phi }\mathrm{\Gamma }_\mu [a_\mu ,\phi ].$$
(A.18)
Their contribution is
$`{\displaystyle \frac{1}{2}}<V_fV_f>`$ $`=`$ $`64F_2+(16\stackrel{~}{g}_1^{}+8g_2+16F_3+32H_4)`$ (A.19)
$`32F_3+64g_1^{}+32\stackrel{~}{g}_1^{}16g_2+64H_2+64H_3.`$
After summing up (A.10), (A.12), (A.16) and (A.19), we find the 2-loop effective action:
$`\mathrm{\Gamma }_{2loop}`$ $`=`$ $`4F_3+32H_1+32H_2+48H_3+(12+32)H_44H_5`$ (A.20)
$`=`$ $`4F_3.`$
It is because
$`H_1+H_2`$ $`=`$ $`0,`$
$`H_3+H_4`$ $`=`$ $`0,`$
$`H_3H_5`$ $`=`$ $`0.`$ (A.21)
The explicit form of $`F_3`$ is
$`F_3`$ $`=`$ $`{\displaystyle \underset{j_i,p_i,s_i}{}}`$ (A.28)
$`\times {\displaystyle \frac{(2j_1+1)(2p_1+1)(2s_1+1)\left[f_1^4j_1(j_1+1)+f_2^4p_1(p_1+1)+f_3^4s_1(s_1+1)\right]}{\left[f_1^2j_1(j_1+1)+f_2^2p_1(p_1+1)+f_3^2s_1(s_1+1)\right]^2}}`$
$`\times {\displaystyle \frac{(2j_2+1)(2p_2+1)(2s_2+1)}{\left[f_1^2j_2(j_2+1)+f_2^2p_2(p_2+1)+f_3^2s_2(s_2+1)\right]}}`$
$`\times {\displaystyle \frac{(2j_3+1)(2p_3+1)(2s_3+1)}{\left[f_1^2j_3(j_3+1)+f_2^2p_3(p_3+1)+f_3^2s_3(s_3+1)\right]}}`$
$`\times \left\{\begin{array}{ccc}j_1& j_2& j_3\\ l_1& l_1& l_1\end{array}\right\}^2\left\{\begin{array}{ccc}p_1& p_2& p_3\\ l_2& l_2& l_2\end{array}\right\}^2\left\{\begin{array}{ccc}s_1& s_2& s_3\\ l_3& l_3& l_3\end{array}\right\}^2.`$
In the large $`N`$ limit, we can use the following approximations
$`j_i(j_i+1)`$ $``$ $`j_i^2,`$
$`2j_i+1`$ $``$ $`2j_i,`$
$`{\displaystyle \underset{j_i=1}{\overset{2l_1}{}}}`$ $``$ $`{\displaystyle _0^{2l_1}}๐j_i,`$
$`l_1^3\left\{\begin{array}{ccc}j_1& j_2& j_3\\ l_1& l_1& l_1\end{array}\right\}^2`$ $``$ $`W(j_1^2,j_2^2,j_3^2).`$ (A.31)
We also define new variables as
$`X=j_1^2,Y=j_2^2,Z=j_3^2,`$
$`L=p_1^2,M=p_2^2,N=p_3^2,`$
$`A=s_1^2,B=s_2^2,C=s_3^2.`$ (A.32)
Finally $`F_3`$ assumes the following expression
$`F_3`$ $`=`$ $`{\displaystyle \frac{l_1l_2l_3}{(f_1f_2f_3)^{4/3}}}{\displaystyle _0^4}๐X๐Y๐Z๐L๐M๐N๐A๐B๐C`$ (A.33)
$`\times {\displaystyle \frac{q_1^2r_1X+q_2^2r_2L+q_3^2r_3A}{(q_1r_1X+q_2r_2L+q_3r_3A)^2}}`$
$`\times {\displaystyle \frac{1}{(q_1r_1Y+q_2r_2M+q_3r_3B)(q_1r_1Z+q_2r_2N+q_3r_3C)}}`$
$`\times W(X,Y,Z)W(L,M,N)W(A,B,C),`$
where
$`r_1=\left({\displaystyle \frac{l_1l_1}{l_2l_3}}\right)^{\frac{2}{3}},r_2=\left({\displaystyle \frac{l_2l_2}{l_1l_3}}\right)^{\frac{2}{3}},r_3=\left({\displaystyle \frac{l_3l_3}{l_1l_2}}\right)^{\frac{2}{3}},`$
$`q_1=\left({\displaystyle \frac{f_1f_1}{f_2f_3}}\right)^{\frac{2}{3}},q_2=\left({\displaystyle \frac{f_2f_2}{f_1f_3}}\right)^{\frac{2}{3}},q_3=\left({\displaystyle \frac{f_3f_3}{f_1f_2}}\right)^{\frac{2}{3}}.`$ (A.34)
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# 1 Introduction
## 1 Introduction
In very simplified terms, studies of classical statistical systems involves two main domains: (1) investigation of the ground state, and (2) summation over fluctuations around this ground state. Although formally according to the definition of the partition function, one has to perform summation over the whole configurational space of a system, in reality it is never done. And it is not that we are doing something wrong. The point is that in most of the cases only very limited part of the configurational space which is relevant for observable thermodynamics. Very often, it the question, what this โrelevant partโ is (which involves the choice of the so called โrelevant variablesโ), which is the most difficult. Studies of the systems containing quenched disorder, in addition to the two items mentioned above, involves the third one (although, technically, very often it turns into the starting one), which is the averaging of self-averaging quantities over random parameters. Nevertheless, at a qualitative level, the situation here remains the same: only very limited part of the configurational space is relevant for observable thermodynamics.
However, in some statistical systems, besides the ground state another local minimum (or minima) could exists. Let us consider extremely simplified situation, schematically shown in Fig.1, when in addition to the ground state the system has another local minimum located in the configurational space โfar awayโ from the ground states, and separated from it by a big (compared to the temperature) energy barrier, which, however, remains finite in the thermodynamic limit. If we are dealing with the system which contains no quenched disorder, then the thermodynamic contribution due to this another state with an exponential accuracy will be simply of order of $`\mathrm{exp}(\beta \mathrm{\Delta }E)`$ (provided $`\mathrm{\Delta }ET`$), where $`\mathrm{\Delta }E=E_1E_0`$. The crucial point, however, is that to get the above exponential contribution, one has to know about existence of the other local minimum, otherwise, its contribution would be just missing. In other words, summing up the perturbation theory around the ground state, and even taking into account all non-linear terms of the Hamiltonian, responsible for the existence of the other local minimum, would not recover the contribution $`\mathrm{exp}(\beta \mathrm{\Delta }E)`$ of the other state, which is located โbeyond barrierโ. It is these type of contributions which are usually called โoff-perturbativeโ.
In the studies of the effects produced by the quenched disorder, conditionally, one can distinguish two main domains of research: strongly disordered systems, like spin-glasses (where the disorder in the dominant factor), and the systems containing some kind of weak disorder which is supposed not to destroy the ground states properties of the corresponding pure system. Traditionally, magnetic statistical systems containing weak disorder, such as random bond ferromagnetic Ising models, are studied focusing mainly on modifications introduced into their critical properties at the phase transition point. In fact, as was pointed out by Griffith many years ago, modification of the critical behavior, is not the only qualitative physical phenomenon which can be produced here.
Let us come back to the example shown in Fig.1. The presence of weak quenched disorder here, provided it does not ruin the global structure of the phase space, would just require supplementary averaging of the above exponential factor, $`\mathrm{exp}(\beta \mathrm{\Delta }E)`$ (since both the energy of the ground state $`E_0`$ and the energy of the excited state $`E_1`$ are now the functions of the disorder parameters), but qualitatively, it would not modify the situation too much. Completely different and new physical phenomena comes into play when the structure of the phase space similar to that shown in Fig.1 is created by the presence of randomness. In other words, this is the situation when weak quenched disorder, although it does not modify the ground state of the system, creates something completely new (absent in the corresponding pure system), namely, the local minima states, somewhere at the periphery of the phase space, โfar awayโ from the ground state of the system.
According to the original observation by Griffiths and later studies, the presence of such off-perturbative states in the disordered ferromagnetic Ising model makes its free energy to be non-analytic function of the external magnetic field $`h`$ in a whole temperature interval above the ferromagnetic phase transition point. Moreover, at least in some cases, this non-analyticity has the form of essential singularity in the limit $`h0`$, . It has to be stressed that all such contributions are just missing in the traditional (perturbative) RG treatment of the problem. In more spectacular way the off-perturbative effects manifest themselves in the dynamical properties, producing the slowering down of the relaxation processes, as well as in the quantum systems (see e.g. and references therein).
Although at a qualitative level the origin of the off-perturbative contributions is more or less clear, their technical implementation, namely, the derivation of e.g. the non-analytic part of the free energy, turns out to be rather tricky problem. Usually, analytic calculations in disordered systems are performed in terms of the replica method. Many years ago Parisi has suggested, that the presence of additional local minima configurations in weakly disordered systems is related, in the replica approach, to the existence of localized in space and breaking replica symmetry instanton-like excitations (translation invariance and replica symmetry is recovered by taking into account all possible excitations of this kind). Later on there were several attempts of concrete realization of this idea for the random temperature and the random field Ising models where it has been demonstrated that the corresponding saddle-point equations may indeed have instanton-like solutions. Next step has been done when the systematic method of summation over all such type of solutions, breaking symmetry in the replica vector order parameter, has been developed. In terms of this method the explicit form of the off-perturbative contributions in the random temperature Ising model has been derived.
In the present paper (following the recent study the original Griffith problem) the systematic approach for the off-perturbative calculations and its relation with the method of the vector replica symmetry breaking is formulated (section II). Then, proposed scheme is applied for the random temperature (section III) and the random field (section IV) ferromagnetic Ising models. It is shown that in both systems at temperatures away from the critical point, namely, in the paramagnetic phase of the random temperature model, and in the ferromagnetic phase of the random field one, the free energy contains non-analytic contributions which (as the functions of the disorder parameters) have the form of essential singularities. It is demonstrated that these contributions appear due to localized in space instanton-like configurations. Physical discussion of the obtained results is given in Section V.
## 2 General scheme of calculations
Let us consider a general (continuous) $`D`$-dimensional random system described by a Hamiltonian $`H[\varphi (๐ฑ);\xi (๐ฑ)]`$, where $`\varphi (๐ฑ)`$ is a field which defines the microscopic state of the system, and $`\xi (๐ฑ)`$ are quenched random parameters. Let us suppose that in addition to the ground state, this system has another thermodynamically relevant (Griffith) region of the configurational space located โfar awayโ from the ground state and separated from it by a finite barrier of the free energy (see Fig.1). In other words, it is supposed that the partition function (of a given sample) can be represented in the form of two separate contributions:
$$Z=๐\varphi (๐ฑ)e^{\beta H}=e^{\beta F_0}+e^{\beta F_1}Z_0+Z_1$$
(1)
where $`F_0`$ is the contribution coming from the vicinity of the ground state, and $`F_1`$ is the contribution of the Griffiths region. Then, for the averaged over disorder total free energy we find:
$$=\frac{1}{\beta }\overline{\mathrm{ln}Z}=\overline{F_0}\frac{1}{\beta }\overline{\mathrm{ln}\left[1+Z_1Z_0^1\right]}$$
(2)
The second term in the above equation, which will be denoted by $`F_G`$, can be represented in the form of the series:
$$F_G=\frac{1}{\beta }\underset{m=1}{\overset{\mathrm{}}{}}\frac{(1)^{m1}}{m}\overline{Z_1^mZ_0^m}=\frac{1}{\beta }\underset{n0}{lim}\underset{m=1}{\overset{\mathrm{}}{}}\frac{(1)^{m1}}{m}Z_n(m)$$
(3)
where
$$Z_n(m)=\underset{b=1}{\overset{m}{}}๐\varphi _b^{(1)}\underset{c=1}{\overset{nm}{}}๐\varphi _c^{(0)}e^{\beta H_n[\varphi _1^{(1)},\mathrm{},\varphi _m^{(1)},\varphi _1^{(0)},\mathrm{},\varphi _{nm}^{(0)}]}$$
(4)
is the replica partition function ($`H_n\left[\mathit{\varphi }\right]`$ is the corresponding replica Hamiltonian), in which the replica symmetry in the $`n`$-component vector field $`\varphi _a`$ ($`a=1,\mathrm{},n`$) is assumed to be broken. Namely, it is supposed that the saddle-point equations
$$\frac{\delta H_n\left[\mathit{\varphi }\right]}{\delta \varphi _a(๐ฑ)}=\mathrm{\hspace{0.33em}0},(a=1,\mathrm{},n)$$
(5)
have non-trivial solutions with the RSB structure
$$\varphi _a^{}(๐ฑ)=\{\begin{array}{cc}\varphi _1(๐ฑ)\hfill & \text{for }a=1,\mathrm{},m\hfill \\ & \\ \varphi _0(๐ฑ)\hfill & \text{for }a=m+1,\mathrm{},n\hfill \end{array}$$
(6)
with $`\varphi _1(๐ฑ)\varphi _0(๐ฑ)`$, so that the integration in the above partition function, eq.(4), goes over fluctuations in the vicinity of these components:
$`\varphi _b^{(1)}(๐ฑ)`$ $`=`$ $`\varphi _1(๐ฑ)+\phi _b(๐ฑ),(b=1,\mathrm{},m)`$
$`\varphi _c^{(0)}(๐ฑ)`$ $`=`$ $`\varphi _0(๐ฑ)+\chi _c(๐ฑ),(c=1,\mathrm{},nm)`$ (7)
It should be stressed that to be thermodynamically relevant, the RSB saddle-point solutions, eq.(6), should satisfy the following tree crucial conditions:
(1) the solutions should be localized in space, so that they are characterized by finite space sizes $`R(m)`$; in this case the partition function, eq.(4), will be proportional to the entropy factor $`V/R^D(m)`$ (where $`V`$ is the volume of the system), and the corresponding free energy contribution $`F_G`$, eq.(3), will be extensive quantity;
(2) they should have finite energies $`E(m)=H_n\left[\mathit{\varphi }^{}\right]`$;
(3) the corresponding Hessian matrix of these solutions should have all eigenvalues positive.
Thus, in the systematic calculations one should find all saddle-point RSB solutions $`\varphi _a^{}(๐ฑ)`$, eq.(6), (satisfying the above three requirements), after that one has to compute their energies $`E(m)`$ (for $`n0`$), next one has to integrate over the fluctuations in the vicinity of these solutions, and finally one has to sum up the series
$$F_G=\frac{V}{\beta }\underset{m=1}{\overset{\mathrm{}}{}}\frac{(1)^{m1}}{m}R^D(m)\left(\beta det\widehat{T}\right)_{n=0}^{1/2}e^{\beta E(m)}$$
(8)
where $`\widehat{T}`$ is the ($`n\times n`$) matrix
$$T_{aa^{}}=\frac{\delta ^2H\left[\mathit{\varphi }\right]}{\delta \varphi _a\delta \varphi _a^{}}|_{\mathit{\varphi }=\mathit{\varphi }^{}}$$
(9)
Note that in the present approach the procedure of analytic continuation $`n0`$ is quite similar to that in the usual replica theory : whenever the parameter $`n`$ becomes an algebraic factor (and not the summation parameter, or the matrix size, etc.), it can safely be set to zero right away.
The above scheme of calculations can be easily generalized for an arbitrary number of the Griffiths regions. For example, let us consider the situation which is qualitatively represented in Fig.2, when in addition to the ground state, the system has two thermodynamically relevant Griffiths states. In this case instead of eq.(1) we will have
$$Z=๐\varphi (x)e^{\beta H}=e^{\beta F_0}+e^{\beta F_1}+e^{\beta F_2}Z_0+Z_1+Z_2$$
(10)
and correspondingly, instead of eq.(3) we find
$`F_G`$ $`=`$ $`{\displaystyle \frac{1}{\beta }}\overline{\mathrm{ln}\left[1+Z_1Z_0^1+Z_2Z_0^1\right]}={\displaystyle \frac{1}{\beta }}{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^{m1}}{m}}{\displaystyle \underset{k=0}{\overset{m}{}}}{\displaystyle \frac{m!}{k!(mk)!}}\overline{\left(Z_1^kZ_2^{mk}Z_0^m\right)}`$ (11)
$`=`$ $`{\displaystyle \frac{1}{\beta }}\underset{n0}{lim}{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^{m1}}{m}}{\displaystyle \underset{k=0}{\overset{m}{}}}{\displaystyle \frac{m!}{k!(mk)!}}Z_n(k,m)`$
Here, in the replica partition function
$$Z_n(k,m)=\underset{b=1}{\overset{k}{}}๐\varphi _b^{(1)}\underset{c=1}{\overset{mk}{}}๐\varphi _c^{(2)}\underset{d=1}{\overset{nm}{}}๐\varphi ^{(0)}e^{\beta H_n[\mathit{\varphi }^{(1)},\mathit{\varphi }^{(2)},\mathit{\varphi }^{(0)}]}$$
(12)
the integration is supposed to be performed in the vicinity of the saddle-point replica vector
$$\varphi _a^{}(๐ฑ)=\{\begin{array}{cc}\varphi _1(๐ฑ),\hfill & \text{for }a=1,\mathrm{},k\hfill \\ & \\ \varphi _2(๐ฑ),\hfill & \text{for }a=k+1,\mathrm{},m\hfill \\ & \\ \varphi _0(๐ฑ),\hfill & \text{for }a=m+1,\mathrm{},n\hfill \end{array}$$
(13)
(where $`\varphi _1(๐ฑ)\varphi _2(๐ฑ)\varphi _0(๐ฑ)`$) which is the solution of the saddle-point equations (5). Finally, for the Griffiths free energy contribution, instead of eq.(8) one obtain
$$F_G=V\underset{m=1}{\overset{\mathrm{}}{}}\frac{(1)^{m1}}{\beta m}\underset{k=0}{\overset{m}{}}\frac{m!}{k!(mk)!}R^D\left(\beta det\widehat{T}\right)_{n=0}^{1/2}e^{\beta E(k,m)}$$
(14)
where $`E(k,m)=H_{n0}\left[\mathit{\varphi }^{}\right]`$ is the energy of a given solution, eq.(13), and $`\widehat{T}`$ is the Hessian matrix, eq.(9).
It is interesting to note that one can arrive to the same representations for the off-perturbative free energy contributions, eq.(8) or eq.(14), following the so called vector replica symmetry breaking scheme. The starting point here is the standard replica definition for the averaged over disorder free energy,
$$=\frac{1}{\beta }\underset{n0}{lim}\frac{\overline{Z^n}1}{n}$$
(15)
where the replica partition function $`\overline{Z^n}Z_n`$ is formally defined by the integration over all configurational space:
$$Z_n=\underset{a=1}{\overset{n}{}}๐\varphi _ae^{\beta H_n[\varphi _1,\mathrm{},\varphi _n]}$$
(16)
Now, let us suppose that in addition to the usual replica symmetric (RS) ground state configuration, the saddle-point equations (5) have another types of solutions, which are well separated in the configurational space from the RS state. In this case (again, denoting their contributions by the label โGโ) the replica partition function, eq.(16), can be decomposed into two parts:
$$Z_n=Z_{RS}+Z_G$$
(17)
Here $`Z_{RS}`$ contains all โroutineโ perturbative contributions in the vicinity of the ground state, and, as usual, (in the limit $`n0`$) this partition function can be represented in the form:
$$Z_{RS}=e^{\beta nF_0}$$
(18)
Thus, according to eq.(15) for the total free energy we get:
$$=F_0+F_G$$
(19)
where
$$F_G=\underset{n0}{lim}\frac{1}{\beta n}Z_G$$
(20)
contains all non-replica-symmetric contributions (if any). As an example, let us suppose that the saddle-point eqs.(5) have non-trivial solutions with the three groups structure like in eq.(13). Moreover, let us suppose that these solutions possess three crucial properties: (1) they are localized in space and characterized by finite spatial sizes $`R_n(m,k)`$; (2) they have finite energies $`E_n(k,m)`$; and (3) their Hessian matrices $`T_{ab}`$, eq.(9), have all the eigenvalues positive. Than taking into account all possible permutations of the tree replica vector components the above free energy $`F_G`$ can be represented in the form
$$F_G=\underset{n0}{lim}\frac{1}{\beta n}\underset{m=1}{\overset{n}{}}\frac{n!}{m!(nm)!}\underset{k=0}{\overset{m}{}}\frac{m!}{k!(mk)!}\frac{V}{R_n^D(k,m)}\left(\beta det\widehat{T}\right)^{1/2}e^{\beta E_n(k,m)}$$
(21)
To perform the analytic continuation $`n0`$ in the above expression the parameter $`n`$ must enter as an algebraic factor, and not as the parameter of summation. This can be achieved if we represent the above series in the following way:
$$F_G=\underset{n0}{lim}\frac{1}{\beta n}\underset{m=1}{\overset{\mathrm{}}{}}\frac{\mathrm{\Gamma }(n+1)}{\mathrm{\Gamma }(m+1)\mathrm{\Gamma }(nm+1)}\underset{k=0}{\overset{m}{}}\frac{m!}{k!(mk)!}\frac{V}{R_n^D(k,m)}\left(\beta det\widehat{T}\right)^{1/2}e^{\beta E_n(k,m)}$$
(22)
Here the summation over $`m`$ is extended beyond $`m=n`$ limit since the gamma function $`\mathrm{\Gamma }(z)`$is equal to infinity both at $`z=0`$ and at all negative integers. Now using the relation:
$$\mathrm{\Gamma }(z)=\frac{\pi }{z\mathrm{\Gamma }(z)\mathrm{sin}(\pi z)}$$
(23)
and referring to โgoodโ analytical properties of the Gamma functions, we can perform the analytic continuation $`n0`$:
$`{\displaystyle \frac{\mathrm{\Gamma }(n+1)}{\mathrm{\Gamma }(m+1)\mathrm{\Gamma }(nm+1)}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(n+1)}{\mathrm{\Gamma }(m+1)\mathrm{\Gamma }[(m1n)]}}`$ (24)
$`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(n+1)(m1n)\mathrm{\Gamma }(m1n)\mathrm{sin}[\pi (m1)\pi n]}{\pi \mathrm{\Gamma }(m+1)}}`$
$`=`$ $`{\displaystyle \frac{\mathrm{sin}(\pi n)}{\pi }}(1)^{m1}{\displaystyle \frac{\mathrm{\Gamma }(n+1)\mathrm{\Gamma }(mn)}{\mathrm{\Gamma }(m+1)}}|_{n0}`$
$``$ $`n{\displaystyle \frac{(1)^{m1}}{m}}`$
Substituting this into eq.(22) we obtain eq.(14) (where $`R(k,m)R_{n=0}(k,m)`$ and $`E(k,m)E_{n=0}(k,m)`$)
Now, in the next two sections I am going to demonstrate how the above general scheme works in the concrete cases of the random temperature and the random field ferromagnetic Ising models.
## 3 Random temperature Ising model
Let us consider weakly disordered $`D`$-dimensional Ising model described by the continuous Ginzburg-Landau Hamiltonian:
$$H=d^Dx\left[\frac{1}{2}\left(\varphi (๐ฑ)\right)^2+\frac{1}{2}\left(\tau \delta \tau (๐ฑ)\right)\varphi ^2(๐ฑ)+\frac{1}{4}g\varphi ^4(๐ฑ)\right]$$
(25)
The disorder is modeled here by a random function $`\tau (๐ฑ)`$ which is described by the Gaussian distribution,
$$P[\delta \tau ]=p_0\mathrm{exp}\left(\frac{1}{4u}d^Dx(\delta \tau (๐ฑ))^2\right),$$
(26)
where $`u`$ is the small parameter which describes the strength of the disorder, and $`p_0`$ is the normalization constant. We are going to consider this system in the paramagnetic phase away from the critical point, so that the reduced temperature parameter $`\tau `$ will be taken to be positive and not too small (it will be demonstrated below that in dimensions $`D<4`$ the limitation on the value of $`\tau `$ is given by the usual Ginzburg-Landau condition, $`\tau g^{2/(4D)}`$).
Weakly disordered systems described in terms of the continuous Ginzburg-Landau Hamiltonian have been usually studied in the framework of the renormalization-group (RG), (perturbative) treatment. In this approach one is able to perform the systematic integration over all fluctuations at the background of the homogeneous state up to the scales of the correlation length $`R_c(\tau )`$. In some cases, this makes possible to derive the leading singularities of the thermodynamical functions in the critical point, at $`\tau 0`$, where $`R_c(\tau )`$ diverges.
However, considering paramagnetic phase of this system, intuitively, it is clear that the contributions of non-homogeneous local minima configurations at scales bigger that the correlation length, which exist due to rare localized in space โferromagnetic islandsโ with negative effective value of the local temperature parameter $`(\tau \delta \tau )`$ are missing in the traditional RG treatment of the problem. To what extend these states are relevant for the critical behavior (at $`\tau 0`$) is still unclear. In the present study , however, I am going to address much more simple question: what is the explicit form of the free energy contributions due to such off-perturbative states away from $`T_c`$, (at $`\tau \tau _g`$).
In fact, at purely heuristic level, it is not so difficult to estimate form of these contributions. Let us consider the spatial island of the linear size $`L`$ characterized by the typical value of the โlocal temperatureโ $`(\tau \delta \tau )=\xi <0`$ Its probability is exponentially small,
$$๐ซ[L,\xi ]\mathrm{exp}\left(\frac{(\tau +\xi )^2}{4u}L^D\right),$$
(27)
and therefore such islands are well separated from each other and can be considered as non-interacting.
It has to be noted that the island with small (negative) value of the local temperature parameter $`\xi `$ can be characterized by the mean-field โupโ and โdownโ states only if its size is bigger than its local correlation length $`R_c(\xi )\xi ^{1/2}`$. Thus, the contribution to the free energy coming from the local ferromagnetic states of such islands with the exponential accuracy can be estimated by their probability:
$$F_G_0^{\mathrm{}}๐\xi _{R_c(\xi )}^{\mathrm{}}๐L\mathrm{exp}\left[\frac{1}{4u}(\tau +\xi )^2L^D\right]_0^{\mathrm{}}๐\xi \mathrm{exp}\left[(const)\frac{(\tau +\xi )^2}{u}\xi ^{D/2}\right]$$
(28)
Here in the integration over $`\xi `$ the leading contribution comes from the vicinity of the saddle-point value
$$\xi _{}=\frac{D}{4D}\tau $$
(29)
(which is positive in dimensions $`D<4`$, and $`\xi _{}\tau _g`$ provided $`\tau \tau _g`$). In this way we obtain the following estimate for the off-perturbative contributions coming from rare locally ferromagnetic islands:
$$F_G\mathrm{exp}\left[(const)\frac{\tau ^{(4D)/2}}{u}\right]$$
(30)
Now let us consider how this result (including the value of the $`(const)`$ factor) can be derived analytically in terms of the systematic approach developed in the previous section. This derivation has been already reported elsewhere. Here I am going to give some more details about the corresponding replica instanton solutions, but in general this section can be considered just as a โwarming upโ exercise before passing to more difficult calculations for the random field model considered in the next section.
Performing the standard Gaussian integration over random parameters $`\delta \tau (๐ฑ)`$, for the replica partition function one gets
$$Z_n=\underset{a=1}{\overset{n}{}}๐\varphi _a(x)e^{\beta H_n\left[\mathit{\varphi }\right]}$$
(31)
where
$$H_n\left[\mathit{\varphi }\right]=d^Dx\left[\frac{1}{2}\underset{a=1}{\overset{n}{}}\left(\varphi _a\right)^2+\frac{1}{2}\tau \underset{a=1}{\overset{n}{}}\varphi _a^2+\frac{1}{4}g\underset{a=1}{\overset{n}{}}\varphi _a^4\frac{1}{4}u\underset{a,b=1}{\overset{n}{}}\varphi _a^2\varphi _b^2\right]$$
(32)
is the corresponding replica Hamiltonian. The saddle-point configurations of the fields $`\varphi _a(๐ฑ)`$ are defined by the equations
$$\mathrm{\Delta }\varphi _a(๐ฑ)+\tau \varphi _a(๐ฑ)+g\varphi _a^3(๐ฑ)u\varphi _a(๐ฑ)\left(\underset{b=1}{\overset{n}{}}\varphi _b^2(๐ฑ)\right)=0$$
(33)
Below we are going to demonstrate that besides the trivial solution $`\varphi _a(๐ฑ)=0`$ these equations have non-trivial localized in space instanton-like solutions with the RSB (two groups) structure:
$$\varphi _a^{}(๐ฑ)=\{\begin{array}{cc}\varphi _1(๐ฑ)\hfill & \text{for }a=1,\mathrm{},m\hfill \\ & \\ 0\hfill & \text{for }a=m+1,\mathrm{},n\hfill \end{array}$$
(34)
Substituting this anzats into the saddle-point eqs(33) and into the Hamiltonian, eq.(32), we find that (in the limit $`n0`$) the instanton configuration $`\varphi _1(x)`$ is defined by the equation
$$\mathrm{\Delta }\varphi _1(๐ฑ)+\tau \varphi _1(๐ฑ)\lambda (m)\varphi _1(๐ฑ)^3=0$$
(35)
which is controlled by the parameter
$$\lambda (m)=umg$$
(36)
and the energy of this configuration is
$$E(m)=md^Dx\left[\frac{1}{2}(\varphi _1)^2+\frac{1}{2}\tau \varphi _1^2\frac{1}{4}\lambda (m)\varphi _1^4\right]$$
(37)
In what follows the parameter $`\lambda (m)`$ will be assumed to be positive. In other words, the solution, which we are going to derived below, exists only for $`m`$ such that $`m>[g/u]`$ (where $`[\mathrm{}]`$ denotes the integer part). It has to be noted that one should not be confused by the โwrongโ sign of the coupling $`\varphi ^4`$ term in the above equations. In fact, it can be shown that the integration over the replica fluctuations around considered solution in the limit $`n0`$ yields the Hessian matrix which has all the eigenvalues positive (this is the usual situation for the replica theory, where the minima of the physical quantities in the limit $`n0`$ turns into maxima of the corresponding replica quantities).
Rescaling the fields,
$$\varphi _1(๐ฑ)=\sqrt{\frac{\tau }{\lambda (m)}}\psi (๐ฑ\sqrt{\tau })$$
(38)
and introducing $`๐ณ๐ฑ\sqrt{\tau }`$, instead of eq.(35) one get the differential equation which contains no parameters:
$$\mathrm{\Delta }\psi (๐ณ)+\psi (๐ณ)\psi ^3(๐ณ)=0$$
(39)
Correspondingly, for the energy of this configuration, eq.(37), one obtains:
$$E(m)=\frac{m}{umg}\tau ^{(4D)/2}E_0(D)$$
(40)
where the quantity $`E_0(D)`$ depends only on the dimensionality of the system:
$$E_0(D)=d^Dz\left[\frac{1}{2}(\psi (๐ณ))^2+\frac{1}{2}\psi ^2(๐ณ)\frac{1}{4}\psi ^4(๐ณ)\right]$$
(41)
It can be shown (see e.g. ) that in dimensions $`D<4`$ eq.(39) has the smooth (with $`\psi ^{}(0)=0`$) spherically symmetric instanton-like solution $`\psi (r)`$ (where $`r=|๐ณ|`$) such that:
$`\psi (r1)`$ $``$ $`\psi (r=0)\psi _01,`$
$`\psi (r1)`$ $``$ $`e^r0.`$ (42)
The energy $`E_0(D)`$ of this solution is a finite and positive number. In particular, in dimensions $`D=3`$, $`\psi _04.34`$ and $`E_018.9`$ (see Fig.3).
As the dimension parameter $`D`$ approaches the upper critical dimensionality $`D_c=4`$, from below the value of the field at the origin $`\psi _0(D)`$ tend to infinity (see Fig.4), while the energies of the corresponding instanton configurations $`E_0(D)`$ approach the finite universal value $`E_0(D4)=E_{}26.3`$. Above dimensions $`D=4`$ eq.(39) has no smooth instanton-like solutions. In other words, described in terms of the dimensions parameter $`D`$, when passing the critical value $`D_c=4`$ from below, the instanton solution disappears in the discontinuous way.
Note that according to the rescaling, eq.(38), the size of these instanton solutions in terms of the original fields $`\varphi _1(๐ฑ)`$ is $`R_c(\tau )=\tau ^{1/2}`$ (which is the usual correlation length of the Ginsburg-Landau theory) and it does not depends on $`m`$. Note also that due to obvious symmetry property $`\varphi _a\varphi _a`$ of the original saddle-point eqs.(33) (which is valid for all non-zero replica field components independently), the above instanton solution $`\varphi _a^{}(๐ฑ)`$, eq.(34), has additional degeneracy factor $`2^m`$.
The final step is the integration over fluctuations $`\phi _a(๐ฑ)`$ at the background of the above instanton solution. Substituting $`\varphi _a(๐ฑ)=\varphi _a^{}(๐ฑ)+\phi _a(๐ฑ)`$, and expanding the Hamiltonian up to the second order in $`\phi _a(๐ฑ)`$, one has to perform the standard Gaussian integration. These calculations, although slightly cumbersome, are quite straightforward (for the details see Refs.). In the result for the Hessian factors one gets
$$\left(det\widehat{T}\right)_{n0}^{1/2}\mathrm{exp}\left[\frac{3m}{2(umg)}g\psi _o^2\right]$$
(43)
Comparing this with factor $`\mathrm{exp}[E(m)]`$, where $`E(m)`$ is the instanton energy, eq.(40), we see, that under condition
$$\tau \tau _g=g^{2/(4D)}$$
(44)
the contribution of fluctuations can be neglected. This is not surprising because eq.(44) is nothing else, but the Ginzburg-Landau criterion which defines the temperature region away from $`T_c`$, where the critical fluctuations are irrelevant. On the other hand, it has to be stressed that in the close vicinity of $`T_c`$ (at $`\tau \tau _g`$), where the critical fluctuations are relevant, the Gaussian approximation used for obtaining the result, eq.(43), can not be valid anymore, and to derive the corresponding fluctuations contribution one would have to start some kind of RG procedure which would properly take into account non-Gaussian interactions. Thus, the above result for the fluctuations contribution, eq.(43), either can be considered as the small correction (at $`\tau \tau _g`$), or otherwise (if it is not small), it is not valid (at $`\tau \tau _g`$).
Thus, substituting the value of the instanton energy, eq.(40), its size $`R=\tau ^{1/2}`$ as well as its degeneracy factor $`2^m`$ into the series, eq.(8), we get
$$F_GV\tau ^{D/2}\underset{m=[g/u]+1}{\overset{\mathrm{}}{}}\frac{(1)^{m1}}{m}\mathrm{\hspace{0.33em}2}^m\mathrm{exp}\left[E_0(D)\frac{m}{umg}\tau ^{(4D)/2}\right]$$
(45)
The exact summation of this series seems to be rather tricky problem, but with the exponential accuracy it can be estimated in a very simple way. One can easily see that in the limit of weak disorder, at $`ug`$, the leading contribution in this summation comes from the region $`mg/u1`$ (where the exponential factor in eq.(45) becomes $`m`$-independent) and this contribution is
$$F_G\mathrm{exp}\left(E_0(D)\frac{\tau ^{(4D)/2}}{u}\right)$$
(46)
We see that obtained off-perturbative (Griffith-like) part of the free energy, as the function of the disorder parameter in the limit $`u0`$, has the form of the essential singularity. Note again that this contribution exists only in dimensions $`D<4`$. As discussed above, at $`D4`$ (the upper critical dimensions), the dimensionless instanton energy factor $`E_0(D)`$ approaches the finite universal limiting value $`E_{}26.3`$. In our world, in three dimensions, $`E_0(D=3)18.9`$.
## 4 Random field Ising model
To study the off-perturbative effects in the random field Ising model we are going to use again the Ginzburg-Landau continuous representation:
$$H=d^Dx\left[\frac{1}{2}\left(\varphi (๐ฑ)\right)^2+\frac{1}{2}\tau \varphi ^2(๐ฑ)+\frac{1}{4}g\varphi ^4(๐ฑ)h(x)\varphi (๐ฑ)\right]$$
(47)
Here the random function $`h(๐ฑ)`$ is described by the Gaussian distribution,
$$P[h(๐ฑ)]=p_0\mathrm{exp}\left(\frac{1}{2h_0}d^Dxh^2(๐ฑ)\right),$$
(48)
where $`h_0`$ is the small parameter which describes the effective strength of the random field, and $`p_0`$ is the normalization constant. Unlike the random temperature model, considered in the previous section, here we are going to consider the system in the low-temperature ferromagnetic phase (supposing that the dimensionality $`D`$ is such that this phase exists), so that the reduced temperature parameter $`\tau `$ will be taken to be negative, $`\tau =|\tau |`$. Again, we will place the system away from the critical point, assuming that the absolute value $`|\tau |`$ is not too small. As usual, to avoid the effects of the critical fluctuations (in dimensions $`D<4`$) we impose the condition $`|\tau |g^{2/(4D)}`$.
Let us suppose that in the absence of the random fields the ferromagnetic ground state of the system, eq.(47), is โupโ. This state (at the mean-field level) is characterized by the order parameter
$$\varphi _0=+\sqrt{\frac{|\tau |}{g}}$$
(49)
and the energy density
$$ฯต_0=\frac{\tau ^2}{4g}$$
(50)
In the usual perturbative approach the effects produced by the random field term of the Hamiltonian, eq.(47), together with the thermal fluctuations could be calculated in the systematic way in terms of the RG procedure (see e.g. and references therein). This approach is designed to take into account all degrees of freedom at scales less that the correlation length, $`R_c(\tau )|\tau |^{1/2}`$.
On the other hand, at scales bigger that the correlation length we can observe completely different type of thermal excitations. Let us consider a spatial island of the linear size $`L`$, where the average value of the field $`h`$ is negative and its absolute value is not too small. Then, in addition to the state โupโ (with slightly modified value of the order parameter), another local minimum with orientation โdownโ can exist in this island (Fig.5). To be stable, the gain in the energy due to the interaction with the field,
$$E_hL^D|h|\varphi _0$$
(51)
should overrun the loss of energy due to the creation of the domain wall,
$$E_{d.w.}L^{D1}\frac{\varphi _0^2}{R_c}$$
(52)
Thus, such double-state situation in the considered island is created provided
$$|h|>\frac{\varphi _0}{LR_c}\frac{|\tau |}{L\sqrt{g}}h_c$$
(53)
According to eq.(48) the probability to find an island of the size $`L`$ with the average value of the field $`h`$ is
$$P(L,h)\mathrm{exp}\left[\frac{h^2}{2h_0^2}L^D\right]$$
(54)
Then the contribution to the free energy of such rare โflippedโ states can be estimated by their probability:
$`F_G`$ $``$ $`{\displaystyle _{R_c}^{\mathrm{}}}๐L{\displaystyle _{h_c}^{\mathrm{}}}๐h\mathrm{exp}\left[{\displaystyle \frac{h^2}{2h_0^2}}L^D\right]`$ (55)
$``$ $`{\displaystyle _{R_c}^{\mathrm{}}}๐L\mathrm{exp}\left[(const){\displaystyle \frac{\tau ^2}{h_0^2g}}L^{D2}\right]`$
$``$ $`\mathrm{exp}\left[(const){\displaystyle \frac{\tau ^{\frac{6D}{2}}}{h_0^2g}}\right]`$
Note that to obtain this result in the above integration over $`L`$ the dimensionality of the system $`D`$ must be bigger than two (otherwise the integral will become divergent). This is nothing else but slightly modified version of the good old Imri-Ma arguments which tells that in dimension $`D2`$ flipping of magnetizations in big spatial islands can become energetically favorable, which indicate the instability of the global ferromagnetic state. Here we assume that the ferromagnetic state is stable, and we see that rare off-perturbative flipping excitations produce non-analytic contribution to the free energy, which in the limit $`h_00`$ has the form of essential singularity.
Now we are going to re-derive the above prediction, eq.(55), in terms of much more rigorous systematic procedure described in section II. Coming back to the original Hamiltonian, eq.(47), after the Gaussian averaging of the replicated partition function over the random function $`h(๐ฑ)`$, one obtains the replica Hamiltonian
$$H_n\left[\mathit{\varphi }\right]=d^Dx\left[\frac{1}{2}\underset{a=1}{\overset{n}{}}\left(\varphi _a\right)^2\frac{1}{2}|\tau |\underset{a=1}{\overset{n}{}}\varphi _a^2+\frac{1}{4}g\underset{a=1}{\overset{n}{}}\varphi _a^4\frac{1}{2}h_0^2\underset{a,b=1}{\overset{n}{}}\varphi _a\varphi _b\right]$$
(56)
The saddle-point configurations of the fields $`\varphi _a(x)`$ are defined by the equations
$$\mathrm{\Delta }\varphi _a(๐ฑ)|\tau |\varphi _a(๐ฑ)+g\varphi _a^3(๐ฑ)h_0^2\left(\underset{b=1}{\overset{n}{}}\varphi _b(๐ฑ)\right)=0$$
(57)
Below we are going to demonstrate that besides the obvious (replica symmetric) ferromagnetic solution $`\varphi _a(๐ฑ)=\sqrt{|\tau |/g}`$ these equations have non-trivial localized in space instanton-like solutions with the RSB two-groups structure:
$$\varphi _a^{}(๐ฑ)=\{\begin{array}{cc}\sqrt{\frac{|\tau |}{g}}\psi _1(๐ฑ\sqrt{|\tau |})\hfill & \text{for }a=1,\mathrm{},m\hfill \\ & \\ \sqrt{\frac{|\tau |}{g}}\psi _0(๐ฑ\sqrt{|\tau |})\hfill & \text{for }a=m+1,\mathrm{},n\hfill \end{array}$$
(58)
Substituting these rescaled fields into the saddle-point eqs(57) and into the Hamiltonian, eq.(56), we find that (in the limit $`n0`$) the instanton configuration $`\{\psi _1(๐ณ),\psi _0(๐ณ)\}`$ (where $`๐ณ=๐ฑ\sqrt{|\tau |}`$) is defined by the two equations
$`\mathrm{\Delta }\psi _1\psi _1+\psi _1^3\lambda (m)\left(\psi _1\psi _0\right)`$ $`=`$ $`0`$
$`\mathrm{\Delta }\psi _0\psi _0+\psi _0^3\lambda (m)\left(\psi _1\psi _0\right)`$ $`=`$ $`0`$ (59)
and its energy is
$$E(m)=m\frac{|\tau |^{2D/2}}{g}d^Dz\left[\frac{1}{2}\left[(\psi _1)^2(\psi _0)^2\right]\frac{1}{2}\left[\psi _1^2\psi _0^2\right]+\frac{1}{4}\left[\psi _1^4\psi _0^4\right]\frac{1}{2}\lambda (m)\left[\psi _1\psi _0\right]^2\right]$$
(60)
where
$$\lambda (m)=\frac{h_0^2m}{|\tau |}$$
(61)
We are looking for the localized in space (spherically symmetric) solutions of the eqs.(4), such that the two functions $`\psi _1(r)`$ and $`\psi _0(r)`$ (where $`r=|๐ณ|`$) are different from each other in a finite region of space, and at large distances they both sufficiently quickly approach the same value $`\psi =1`$, so that the integral in eq.(60) will be converging. Simple analysis of the structure of the โpotential energyโ
$$U(\psi _1,\psi _0)=\frac{1}{2}\left[\psi _1^2\psi _0^2\right]+\frac{1}{4}\left[\psi _1^4\psi _0^4\right]\frac{1}{2}\lambda \left[\psi _1\psi _0\right]^2$$
(62)
shows that until the parameter $`\lambda `$ is small (so that the last coupling term in the above expression is just a small correction), the potential $`U(\psi _1,\psi _0)`$ has 9 saddle-points (in the vicinity of the points $`(0;0),(0;\pm 1),(\pm 1;0),(\pm 1,\pm 1)`$ and $`(\pm 1;1)`$). In this situation the two fields $`\psi _1`$ and $`\psi _0`$ are effectively independent, and no instanton-like solutions described above can exist. When increasing the parameter $`\lambda `$, starting from
$$\lambda >\lambda _c0.23$$
(63)
only 5 saddle points of the potential $`U(\psi _1,\psi _0)`$ remains in the plain $`(\psi _1;\psi _0)`$. They have coordinates: $`(0;0),(\pm 1;\pm 1)`$ and $`(\pm \psi _1^{}(\lambda );\pm \psi _0^{}(\lambda ))`$, where $`0<\psi _{1,0}^{}<1`$ (in particular, $`\psi _1^{}(\lambda _c)0.17`$ and $`\psi _0^{}(\lambda _c)0.90`$). It is crucial that at the points $`(\pm \psi _1^{};\pm \psi _0^{})`$ the potential $`U(\psi _1,\psi _0)`$ has the maxima. It is due to the existence of these maxima that at $`\lambda >\lambda _c`$ the instanton solutions become possible.
Let us consider the limit $`\lambda (m)1`$, or
$$mm_c=\left[\lambda _c\frac{|\tau |}{h_0^2}\right]+1$$
(64)
In this limit, according to eqs.(4), the two fields $`\psi _1`$ and $`\psi _0`$ must be close to each other. Redefining,
$`\psi _1(r)`$ $`=`$ $`\psi (r)+{\displaystyle \frac{1}{\lambda }}\chi (r)`$
$`\psi _0(r)`$ $`=`$ $`\psi (r){\displaystyle \frac{1}{\lambda }}\chi (r)`$ (65)
in the leading order in $`\lambda ^1`$ instead of eqs.(4) we get much more simple equations:
$`\mathrm{\Delta }\psi \psi +\psi ^32\chi `$ $`=`$ $`0`$
$`\mathrm{\Delta }\chi +(3\psi ^21)\chi `$ $`=`$ $`0`$ (66)
which contain no parameters. For the energy of the configurations described by the two fields $`\psi (r)`$ and $`\chi (r)`$ instead of eq.(60) (again, in the leading order in $`\lambda ^1`$) we find the value, which does not depend on the summation parameter $`m`$,
$$E=\frac{|\tau |^{\frac{6D}{2}}}{h_0^2g}E_0(D)$$
(67)
where
$$E_0(D)=d^Dz\left[(\psi )(\chi )+(\psi ^3\psi )\chi \chi ^2\right]$$
(68)
is the universal quantity which depends only on the dimensionality of the system.
It turns out that in dimensions $`D<3`$, the system of eqs.(4), indeed has smooth instanton-like spherically symmetric solution which has finite and positive energy $`E_0(D)`$. Within the limited spatial region $`rr_c1`$, the values of the fields $`\psi (r)`$ and $`\chi (r)`$ are finite and of the order of their values at the origin, $`\psi _01`$ and $`\chi _01`$. On the other hand, at $`rr_c`$ the function $`\psi (r)`$ exponentially quickly approaches 1, while the function $`\chi (r)`$ exponentially tents to zero. The illustration of the instanton solution in the dimension $`D=2.9`$ is given in Fig.6, where $`\psi _00.818`$, $`\chi _00.284`$
Fig.7 demonstrates the corresponding โtrajectoriesโ of the instanton solutions in the plane $`(\psi ,\chi )`$ at various values of the dimension. As the dimension $`D`$ approaches the value $`D_c=3`$ from below, the starting values $`\psi _01`$ and $`\chi _00`$. Above three dimension the instanton solution disappears. Thus we have to conclude that $`D=3`$ is the upper critical dimension for the considered Griffiths phenomena in the RFIM, in agreement with the earlier suggestion as well as with the recent studies of similar instanton-like configurations in the presence of external magnetic field
Let us come back to the general expression for the off-perturbative part of the free energy, the series, eq.(8), where the summation over $`m`$ starts now from $`m=m_c`$, eq.(64). Noting that the instanton energy $`E(m)`$ is the decreasing function of $`m`$, we can conclude that with the exponential accuracy this converging series can be estimated by its asymptotic part at $`mm_c`$. Thus, substituting here the value of the instanton energy, eq.(67), (and neglecting the critical fluctuations), with the exponential accuracy we get the result
$$F_G\mathrm{exp}\left[E_0(D)\frac{|\tau |^{\frac{6D}{2}}}{h_0^2g}\right]$$
(69)
which perfectly agrees with the โhand-wavingโ estimate, eq.(55). This non-analytic in $`h_0`$ contribution, which has the form of essential singularity at $`h_00`$, is valid only in dimensions $`D<3`$, and at temperatures not too close to the critical point, at $`|\tau |g^{2/(4D)}`$. The dimension $`D=3`$ is marginal for this kind of phenomena. Therefore the investigation of the Griffith-like contributions in the three-dimensional RFIM requires much more discreet analysis.
## 5 Discussion
In this paper the systematic method for computing off-perturbative thermodynamic contributions in disordered systems has been proposed. It has been tested on two the most popular classical statistical systems containing quenched disorder: the random temperature and the random field Ising models. In both cases the off-perturbative contributions as the functions of the parameters which describe the strength of the disorder have the form of the essential singularities, eqs.(46), (69).
Of course, thinking about possible experimental or numerical tests, the validity of the obtained results is rather limited. On one hand, since the consideration has been done in terms of the continuous Ginsburg-Landau Hamiltonian, one has to place the system sufficiently close to the phase transition point, so that the correlation length would be large compared to the lattice spacing. On the other hand, since the present study completely neglects the critical fluctuations, the system has to be sufficiently far from the critical point. Formally, in terms of the Ginsburg-Landau Hamiltonian, these two requirements can be easily satisfied: it is sufficient to demand that (1) the coupling parameter $`g`$ is small, and (2) the reduced temperature parameter $`\tau `$ is bounded by the condition $`g^{2/(4D)}|\tau |1`$ (where $`D`$ is the system dimensionality, $`D<4`$). Besides, the strength of the disorder must be small: $`ug`$ in the random temperature model, and $`h_0\sqrt{|\tau |}`$ in the random field one.
On the other hand, keeping in mind more general perspectives, the proposed approach, in my view, may open the way to study the nature of the phase transitions in the considered systems. It is generally believed that it is the off-perturbative states which makes the study of the phase transition in the random field Ising model so difficult. It is remarkable, that according to the present study, quite similar off-perturbative contributions are also present in the random temperature Ising systems where, at least at the qualitative level, the nature of the phase transition was traditionally believed to be well understood (see e.g. ). This supports recent suspects that the off-perturbative effects could be quite relevant for the critical properties of the disordered ferromagnetic systems.
Another interesting observation is that in terms of the considered off-perturbative contributions, the situation, when approaching $`T_c`$ from above and from below, looks totally asymmetric both in the random field and in the random temperature systems. The considered effects are present in the ferromagnetic phase, at $`T<T_c`$, of the random field model, while they are absent in its paramagnetic phase at $`T>T_c`$. On the other hand, the situation in the random temperature model is similar, although โreversedโ: the off-perturbative contributions are present in the paramagnetic phase, at $`T>T_c`$, while they are absent in the ferromagnetic phase, at $`T<T_c`$. May be it is this asymmetry, which makes the nature of the phase transitions in these systems to be so non-trivial?
One more qualitative observation is that according to the present study, the off-perturbative contributions are absent in the random temperature model at dimensions $`D>4`$, and in the random field model at dimensions $`D>3`$. As for the random temperature systems, this is not surprising: all the previous studies were definite that at $`D>4`$ the disorder is irrelevant for the phase transition. What does this mean for the random field model is much less clear, because here it is well established that its upper critical dimensionality is equal to 6 (the dimensionality above which the critical behavior is described by the Gaussian theory, and the presence of the random fields is irrelevant for the phase transition). Well, of course, the absence of the off-perturbative contributions in dimensions $`3<D6`$ does not mean, that the random fields are irrelevant. Probably it indicates that here all the random field effects can be taken into account in the framework of the perturbative RG procedures.
To answer all the above questions (as well as many other important questions which were not formulated here), the only thing which remains to be done is to find the way to overcome the Ginzburg-Landau limitation $`|\tau |g^{2/(4D)}`$, and to take the limit $`\tau 0`$. To do that one has just to formulate a theory which would properly take into account the critical fluctuations on top of the instanton-like configurations described in this paper.
Acknowlegments
The author is grateful to Markus Mueller and Alessandro Silva for their quite helpful comments and advises concerning the numerical solutions for the instanton configurations in the random field Ising model.
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# From Zwiebach invariants to Getzler relation
## 1. Pre-introduction
In , Barannikov and Kontsevich have found a solution to the WDVV equation starting from the algebra of polyvector fields on Calabi-Yau manifolds. The algebraic properties of polyvector fields used in their construction are captured by an abstract algebraic structure called dGBV-algebra with Hodge property.
One of the main results of this paper is a new interpretation of Barannikov-Kontsevich construction. We represent their solution as a sum over trivalent trees. Using this representation we give a new independent proof that this sum over trivalent trees satisfies the WDVV equation.
Since we have a sum over trivalent trees, it is very natural to study the sum over graphs of higher genera (with the same tensor expressions associated to elements of graphs). We prove that in genus $`1`$ our construction satisfied the Getzler elliptic equation . But in order to prove this we have introduced a new surprising algebraic axiom (we call it $`1/12`$ axiom, see Sections 5.3).
Probably, the main problem for us was to find a proper explanation of this additional axiom. In fact, in order to obtain naturally the genus $`0`$ part of our construction (i.e., Barannikov-Kontsevich solution in terms of trivalent trees) it is enough to study the BCOV-action written down in , \[1, Appendix\], , and . But then we have to introduce $`1/12`$-axiom just for computational reasons, and Getzlerโs relation in genus $`1`$ comes over as a miracle.
So, we have found another approach to the explanation of our results. It is a kind of an โoperadicโ homotopy extension of Gromov-Witten theory in the spirit of Getzler and Zwiebach . In this framework all our axioms including $`1/12`$-axiom come very natural. Moreover, in this approach the relations coming from geometry of the moduli space of curves seem to be very expected.
This is an amazing fact that there are two completely different natural approaches to the same construction: one is from the $`B`$-model side (Barannikov-Kontsevich) and another one is from the $`A`$-model side (we call it the theory of Zwiebach invariants).
For introduction we have chosen the second approach, since it better explains our results. However, the origin of the idea to use trivalent trees is also hidden in the first approach and we explain this in the appendix.
## 2. Introduction
String theory appeared in the beginning of seventies as an attempt to find fundamental degrees of freedom that would form theory free of ultraviolet divergences and give gravity as a low energy effective theory.
In its standard formulation the string theory computes $`g`$-loop scattering amplitudes of $`2`$ particles into $`(n2)`$-particles as an integral over the moduli space of complex structures of the genus $`g`$ surface with $`n`$ marked points. The measure of integration is a correlator in a very specific conformal field theory that has an odd symmetry $`Q`$ (such that $`Q^2=0`$) and so-called ghosts (due to gauge fixing of the diffeomorphism invariance). Energy-momentum tensor in such theory is $`Q`$-exact.
In the process of study of string theory it was generalized to the so-called topological string theory. In topological string theory conformal theory with ghosts is replaced by a more general conformal theory with $`Q`$-symmetry and (co)exact energy-momentum tensor.
The most impressive application of these ideas is the theory of geometric Gromov-Witten invariants (known in physics as type $`A`$ topological strings). This theory attracted a lot of attention in last decade since its amplitudes give answers to famous problems in enumerative algebraic geometry.
Further generalization of these ideas involves the construction of the set of factorizable closed forms on the moduli spaces of complex structures on Riemann surfaces (so that the integral of the top form produces amplitudes). In this way we get generalized amplitudes that take values in cohomology of the moduli space. Evaluation of these generalized amplitudes on the contractible cycles (together with the factorization property) leads to relations among amplitudes (WDVV and Getzler relations) that are very important in applications.
This leads to the new definition of amplitudes as a system of factorizable maps from the tensor products of the vector space with bilinear pairing to cohomology of the moduli spaces of Riemann surfaces that we call simply Gromov-Witten invariants. Note, that here we do not insist that amplitudes come from the integral over the moduli space of differential forms coming from conformal field theory; we study amplitudes on their own <sup>1</sup><sup>1</sup>1We call them just Gromov-Witten invariants in order to distinguish them from the geometrical Gromov-Witten invariants that follow from the theory of holomorphic maps..
Note, that formalization of general (irrational) conformal field theory produces objects that are rather dificult to deal with. One has to study infinite sums of tensor products of two irreducible representations of chiral algebras that have to obey additional conditions coming from the duality (see , ). Otherwise, one has to study various limits of rational conformal theories when the number of irreducible representation goes to infinity. It is really a challenge to develop such a theory in full generality and to find a reasonable amount of understandable examples (as far as we know only the theory on a torus and on its orbifolds are constructively known among irrational theories). Such an understanding would be a curious extension of differential geometry but it is out of reach for a moment. Therefore we have to wait a bit before we can say something constructive about general irrational conformal theories with $`Q`$-symmetry and exact energy-momentum tensor.
However, we can say something about the degeneration of this magnificent picture in the limit where the conformal theory degenerates so that the conformal dimensions of some fields tend to zero (note that there are fields with exactly zero dimension among them).
It seems that we can write down tractible axioms on correlators of fields with nearly zero dimension (viewed as differential forms on the proper moduli spaces) at the point of degeneration โ we will call these Zwiebach invariants <sup>2</sup><sup>2</sup>2We call these correlators Zwiebach invariants because of inspiring work of Zwiebach on related issues.. One can show at the heuristic level that the amplitudes in the nearly degenerate theory can be obtained as a sum over graphs with Zwiebach invariants associated to vertices.
Our next step would be to forget about the conformal field theory origin of the procedure and to study the theory of Zwiebach invariants (as a sets of maps taking values in forms on the moduli spaces that obeys some axioms) on their own. It is similar to forgetting the conformal field theory origin of the Gromov-Witten invariants.
However, we will show that now we may also formalize the passage to the nearly degenerate theory, when dimension of some fields is lifted. We will see that this leads to the procedure of induction of the structure of the Zwiebach invariant on the subbicomplexes. And Zwiebach invariants of bicomplex with zero differential turn out to be Gromov-Witten invariants.
After presenting the outline of such general construction we have to study the confirming example โ and we really do it. Namely, we study the case when Zwiebach invariants take the simplest possible form โ they are constructed from the Hodge dGBV algebra that satisfies the $`1/12`$ axiom, and (possibly) some other conditions. Instead of looking for the formal proof that the set of these other conditions is empty (we admit that it would be nice to have such a proof) we just compute directly induced structures on cohomology of the bicomplex in genera zero and one. We show by explicit computation that these structures do solve WDVV and Getzler equations.
All this should be compared with the theory of induction of the homotopical structure on the subcomplex. The simplest homotopical structure is the structure of differential graded (Lie) algebra. Therefore, we propose the generalization of this story to the bicomplexes with dgA being replaced by Hodge dGBV with $`1/12`$ axiom, and homotopical algebra structure being replaced by Zwiebach invariants.
The natural question to ask is whether all Zwiebach invariant can be obtained by induction from the simplest ones (like all homotopical algebras can be obtained by induction from the differential graded ones). We do not know the answer at the moment.
We hope that the notion of Zwiebach invariants will help to understand why constructions of , and lead to Gromov-Witten invariants.
### 2.1. Definition of Gromov-Witten invariants
By Gromov-Witten invariants we mean the set of maps
(1)
$$m_{g,n}:H_0^n๐,$$
where $`H_0`$ is a vector space and $`๐`$ is the space of cycles in the Deligne-Mumford compactification of the moduli space of genus $`g`$ curves with $`n`$ marked points $`\overline{}_{g,n}`$. This set of maps satisfies the following conditions :
1. It is symmetric with respect to diagonal action of the symmetric group on factors of $`H_0^n`$ and cycles in $`\overline{}_{g,n}`$.
2. It vanishes when restricted to cycles that are zero in rational homologies of $`\overline{}_{g,n}`$
3. It satisfies the factorization property described below.
The factorization property corresponds to degenerations of a surface of genus $`g`$ with $`n`$ marked points. First we consider the case when a surface degenerates into surfaces of genera $`g_1`$ and $`g_2`$ with $`n_1`$ and $`n_2`$ marked points respectively and that have a common point:
(2)
$$\begin{array}{c}m_{g,n}(h_1,\mathrm{},h_n)(c_1\times c_2)=\hfill \\ \hfill \underset{i,j}{}\eta ^{ij}m_{g_1,n_1+1}(h_1,\mathrm{},h_{n_1},e_i)(c_1)m_{g_2,n_2+1}(h_{n_1+1},\mathrm{},h_n,e_j)(c_2).\end{array}$$
Here $`\{e_j\}`$ is a basis in $`H_0`$, $`\eta ^{ij}`$ is the inverse metric on $`H_0`$ written in this basis, $`c_1`$ and $`c_2`$ are some cycles in $`\overline{}_{g_1,n_1}`$ and $`\overline{}_{g_2,n_2}`$ respectively, $`c_1\times c_2`$ is viewed as a cycle in $`\overline{}_{g,n}`$ via the embedding $`\overline{}_{g_1,n_1+1}\times \overline{}_{g_2,n_2+1}\overline{}_{g,n}`$.
Then we consider the degeneration of a curve of genus $`g`$ into a curve of genus $`g1`$ with a double point. In this case the factorization property means
(3)
$$m_{g,n}(h_1,\mathrm{},h_n)(c)=\underset{i,j}{}\eta ^{ij}m_{g1,n+2}(h_1,\mathrm{},h_n,e_i,e_j)(c).$$
Here $`c`$ is a cycle in $`\overline{}_{g1,n+2}`$ considered also as a cycle in $`\overline{}_{g,n}`$ via the natural mapping $`\overline{}_{g1,n+2}\overline{}_{g,n}`$.
### 2.2. Set of factorizable maps from topological conformal field theory
In this and in the next subsections we assume some knowledge of conformal field theory (CFT). Reader that does not know CFT may skip this subsection and proceed to Subsection 2.4 where we formalize insights coming from CFT.
Consider CFT with odd symmetry. This means that the space of local observables $`H_c`$ is a complex with the differential $`Q`$, correlators satisfy
(4)
$$\underset{i=1}{\overset{k}{}}v_1,\mathrm{},Q(v_i),\mathrm{},v_k=0,$$
both holomorphic and antiholomorphic energy-momentum tensors are $`Q`$-exact
(5)
$$Q(G)=T,Q(\overline{G})=\overline{T},$$
and the fields $`G`$ and $`\overline{G}`$ do not have singularities in their mutual operator product.
Consider the correlators
(6)
$$v_1(z_1),\mathrm{},v_n(z_n),G(x_1),\mathrm{},G(x_p),\overline{G}(y_1),\mathrm{},\overline{G}(y_q)$$
as differential $`(p,q)`$-forms on the moduli space $`\widehat{}_{g,n}`$ of Riemann surfaces with germs of local coordinates at marked points $`z_1,\mathrm{},z_n`$. This means that we can contract such a form with a holomorphic (and antiholomorphic) vectors, tangent to the moduli space. A holomorphic tangent vector is determined by a Beltrami differential; so we can multiply $`G`$ by the Beltrami differential and integrate over the surface. If $`n`$ is not zero one can also multiply $`G`$ by a holomorphic vector field in the neibourhood of a marked point and integrate around it. Similarly, one can define contraction with an antiholomorphic tangent vector.
This differential form descends down to the moduli space $`_{g,n}`$ if the correlator contains only the first order poles when $`x`$ and $`y`$ approach the set of the marked points. The second order pole in operator product expansion between $`G`$ and $`v`$ is called the action of the operator $`G_0`$ on $`v`$. Similarly, we define $`\overline{G}_0`$.
Only the phase of the local coordinate corresponds to the noncontractable piece of the structure group of the bundle of germs of local coordinates over the moduli space of complex structures with the marked points. Therefore we only have to impose the condition
(7)
$$G_{}(v):=(G_0\overline{G}_0)v=0.$$
In order to get closed forms on the moduli space we impose
(8)
$$Q(v)=0.$$
Finally, we need (and this part of construction is missing in ) our differential form to be extendable to the Deligne-Mumford compactification of the moduli space. One can show that this is satisfied if the fields $`v`$ that are in the image of $`G_{}`$ are not in the kernel of $`T_0+\overline{T}_0`$; here $`T_0=Q(G_0)`$ and $`\overline{T}_0=Q(\overline{G}_0)`$.
### 2.3. Topological string amplitudes in degenerating conformal theory
In the previous subsection we outlined the construction of amplitudes in an arbitrary topological conformal theory. However, in the so called degenerating theories (like type $`B`$ theory on Calaby-Yau at the infinite volume limit) life simplifies a bit, and the construction of amplitudes can be encoded in a tractable linear algebra data.
By a degenerating theory we mean a family of theories parametrized by a parameter $`ฯต`$ such that at $`ฯต=0`$ the subset $`HH_c`$ of fields has zero conformal dimension:
(9)
$$T_0H=\overline{T}_0H=0.$$
In some cases (in particular, in the type $`B`$ example) one can check that in this limit most of the correlators (6) vanish over the bounded domain of moduli of complex structures. However, this does not mean that the integrals over the moduli space vanish. What really happens is the following: the support of the correlation function moves towards the region where surface degenerates. The good model of this phenomena is the ordinary integral:
(10)
$$I(ฯต)=_0^+\mathrm{}\mathrm{exp}(tฯต)ฯต๐t.$$
The value of this integral is independent of $`ฯต`$ while the integrand tends to zero as $`ฯต`$ goes to zero. The support of the integrad is at $`t`$ of order $`\frac{1}{ฯต}`$.
Thus we have a contribution from the boundary of the moduli space (see ). This contribution comes from the infinitely long tubes connecting components of the degenerating surface and equals to
(11)
$$K=\frac{G_0\overline{G}_0}{T_0+\overline{T}_0}=G_{}G_+,$$
where
(12)
$$G_+=\frac{G_0+\overline{G}_0}{T_0+\overline{T}_0}$$
Note, that $`G_+`$ has a regular limit as $`ฯต`$ tends to zero, since
(13)
$$\{Q,G_+\}=1\mathrm{\Pi }_0.$$
where $`\mathrm{\Pi }_0`$ is a projector to the space $`H_0`$ of zero modes of $`T_0+\overline{T}_0`$ that presumably has a smooth limit as $`ฯต`$ goes to zero. Note, that $`H_0`$ is the limit of the kernel rather than the kernel of the limiting operator (that coincides with $`H`$).
Note, that the space $`H`$ is equipped with the bilinear pairing: given by the two-point function:
(14)
$$(v_1,v_2)=v_1(z_1),v_2(z_2).$$
Since the conformal dimension of fields $`v_i`$ is zero, this correlation function is independent of the coordinates $`z_i`$.
Therefore, we obtain the rules for computation of the amplitude in the limiting theory. The contribution from the bulk of the moduli space is obtained by substitution of elements from $`H_0`$. The contribution from the degenerated surfaces is given by the sum over graphs, such that $`k`$-vertices of the graphs are labeled by $`k`$-point correlation functions (of different genera). A weight of a graph is given by the pairing between vertices, so-called propagators $`K`$ (given by (11)) that correspond to edges, and elements of $`H_0`$ that correspond to tails. Pairing is performed with the help of the bilinear form defined in (14). Note, that vertices are paired with $`G_{}`$ closed vectors, therefore vertices correspond to horisontal invariant forms on components of the moduli space and can be integrated over it.
We make an attempt to formalize this in the next subsection.
### 2.4. Zwiebach invariants
In this section, we sketch the principal construction of Zwiebach invariants that motivates our purely algebraic constructions in the rest of the paper. Note, that part of this was already presented in the work of Zwiebach , but he missed the Hodge condition. In different settings, but in a closed way, a piece of the algebraic structure that we finally get was also obtained in .
#### 2.4.1. Kimura-Stasheff-Voronov space
We consider the Kimura-Stasheff-Voronov compactification $`\overline{๐ฆ}_{g,n}`$ of the moduli space of curves of genus $`g`$ with $`n`$ marked point. It is a real blow-up of $`\overline{}_{g,n}`$; we just remember the relative angles at double points. We can also choose an angle of the tangent vector at each marked point; this way we get the principal $`U(1)^n`$-bundle over $`\overline{๐ฆ}_{g,n}`$. We denote the total space of this bundle by $`\overline{๐ฎ}_{g,n}`$.
Let $`H`$ be a bicomplex with two differentials denoted by $`Q`$ and $`G_{}`$ and with a scalar product $`(,)`$ invariant under the differentials: $`(Qv,w)=\pm (v,Qw)`$, $`(G_{}v,w)=\pm (v,G_{}w)`$.
Below we consider the action of $`Q`$ and $`G_{}`$ on $`H^n`$. We denote by $`Q^{(k)}`$ and $`G_{}^{(k)}`$ the action of $`Q`$ and $`G_{}`$ respectively on the $`k`$-th component of the tensor product.
#### 2.4.2. Definition
The Zwiebach invariants is the set $`\{C_{g,n}|g0,n0,3g3+n0\}`$ of $`H^n`$-valued differential forms on $`\overline{๐ฎ}_{g,n}`$, satisfying the axioms:
1. $`C_{g,n}`$ is (graded) symmetric under the interchange of factors in $`H^n`$ with the simultaneous renumeration of marked points;
2. $`C_{g,n}`$ is totally closed, $`(Q+d)C_{g,n}=0`$ ($`Q=_{k=1}^nQ^{(i)}`$);
3. $`C_{g,n}`$ is totally horizontal, $`(G_{}^{(k)}+ฤฑ_k)C_{g,n}=0`$ for all $`1kn`$ (we denote by $`ฤฑ_k`$ the substitution of the vector field generating the action on $`\overline{๐ฎ}_{g,n}`$ of the $`k`$-th copy of $`U(1)`$) and $`C_{g,n}`$ is invariant under the action of $`U(1)^n`$;
4. $`\{C_{g,n}\}`$ is the factorizable set of maps (cf. Equations (2), (3)), that is,
(15) $`C_{g,n}|_{\gamma _2}`$ $`=[C_{g_1,n_1}C_{g_2,n_2}],`$
(16) $`C_{g,n}|_{\gamma _1}`$ $`=[C_{g1,n+2}]`$
Here $`\gamma _2`$ corresponds to the degeneration of the surface into two components, $`\gamma _1`$ corresponds to the degeneration of a handle, and $`[]`$ denotes the contraction with the scalar product of the last factors in $`H^{n_1+1}`$ and $`H^{n_2+1}`$ in the first case and of the last two factors in $`H^{n+2}`$ in the second case.
It is useful to rewrite the last two axioms in local charts. Locally, $`\overline{๐ฎ}_{g,n}`$ is a product of $`\overline{๐ฆ}_{g,n}`$ and $`n`$ circles. Then the horizontality axiom means that $`C_{g,n}`$ is represented as
(17)
$$C_{g,n}=(1+d\varphi _1G_{}^{(1)})\mathrm{}(1+d\varphi _nG_{}^{(n)})\stackrel{~}{C}_{g,n},$$
where $`\stackrel{~}{C}_{g,n}`$ is (the pull-back of) a form on $`\overline{๐ฆ}_{g,n}`$ and $`\varphi _i`$ is the angle at the $`i`$-th marked point. The factorization property in terms of $`\stackrel{~}{C}_{g,n}`$ looks as follows:
(18) $`\stackrel{~}{C}_{g,n}|_{\gamma _2}`$ $`=\left[\stackrel{~}{C}_{g_1,n_1+1}\left(1+d\psi G_{}^{(n_2+1)}\right)\stackrel{~}{C}_{g_2,n_2+1}\right],`$
(19) $`\stackrel{~}{C}_{g,n}|_{\gamma _1}`$ $`=\left[\left(1+d\psi G_{}^{(n+2)}\right)\stackrel{~}{C}_{g1,n+2}\right].`$
Here $`\gamma _2`$ corresponds to the degeneration of the surface into two components, $`\gamma _1`$ corresponds to the degeneration of a handle, $`\psi `$ denotes the relative angle at the double point, and $`[]`$ denotes the contraction with the scalar product of the last factors in $`H^{n_1+1}`$ and $`H^{n_2+1}`$ in the first case and of the last two factors in $`H^{n+2}`$ in the second case. Indeed, we just use that $`\psi =\varphi _{n_1+1}+\varphi _{n_2+1}`$ in the first case and $`\psi =\varphi _{n+1}+\varphi _{n+2}`$ in the second case.
Note that below we usually use $`\stackrel{~}{C}_{g,n}`$ instead of $`C_{g,n}`$ just to make our calculations more transparent.
#### 2.4.3. Gromov-Witten invaiants
Zwiebach invariants on the bicomplex with zero differentials determine Gromov-Witten invariants.
Indeed, in this case $`C_{g,n}=\stackrel{~}{C}_{g,n}`$. Together with the factorization property this means that $`\{C_{g,n}\}`$ is lifted from the the blowdown of Kimura-Stasheff-Voronov spaces, i.e. it is determined by a set of continuous forms on Deligne-Mumford spaces. Therefore, integrating these forms along cycles in $`\overline{}_{g,n}`$ we get Gromov-Witten invariants defined on the space dual to $`H`$.
#### 2.4.4. Induced Zwiebach invariants
Induced Zwiebach invariants are obtained by contraction of an acyclic subbicomplex of $`(H,Q,G_{})`$. Let $`H=H^{}H^{\prime \prime }`$ such that $`(H^{},H^{\prime \prime })=0`$ and $`H^{\prime \prime }`$ is an acyclic subbicomplex. We denote by $`G_+`$ the contraction operator. This means that $`G_+H^{}=0`$, $`\mathrm{\Pi }=\{Q,G_+\}`$ is the projection to $`H^{\prime \prime }`$ along $`H^{}`$, and $`\{G_+,G_{}\}=0`$.
We construct an induced Zwiebach form $`C_{g,n}^{ind}`$ (or rather $`\stackrel{~}{C}_{g,n}^{ind}`$) on a modification of $`\overline{๐ฆ}_{g,n}`$. Each degeneration of a curve gives us a boundary stratum $`\gamma `$ that is a pricipal $`U(1)`$ bundle over $`\overline{๐ฆ}_{g_1,n_1+1}\times \overline{๐ฆ}_{g_1,n_2+1}`$ or $`\overline{๐ฆ}_{g1,n+2}/_2`$ ($`_2`$ exchanges the labels of the last two points). At each such component of the boundary we glue the cylinder $`\gamma \times [0,+\mathrm{}]`$ such that $`\gamma `$ in $`\overline{๐ฆ}_{g,n}`$ is identified with $`\gamma \times \{0\}`$ in the cylinder.
So, we take a form $`\stackrel{~}{C}_{g,n}`$, restrict it to $`H_{}^{}{}_{}{}^{n}`$, and extend it to the cylinder glued at $`\gamma `$ as the restriction to $`H_{}^{}{}_{}{}^{n}`$ of
(20)
$$\left[\stackrel{~}{C}_{g_1,n_1+1}e^{t\mathrm{\Pi }dtG_+}\left(1+d\varphi G_{}^{(n_2+1)}\right)\stackrel{~}{C}_{g_2,n_2+1}\right]$$
in the first case of curve degeneration or
(21)
$$\left[e^{t\mathrm{\Pi }dtG_+}\left(1+d\varphi G_{}^{(n+2)}\right)\stackrel{~}{C}_{g1,n+2}\right]$$
in the second case of curve degeneration. Here $`t`$ is a coordinate along cylinder and operators $`\mathrm{\Pi }`$ and $`G_+`$ in the formulas act at the same copy of $`H`$ as $`G_{}`$. In terms of $`C_{g,n}`$, this is just the same contraction as in Equations (15) and (16), but the scalar product is defined as $`(V,W)_t=(V,e^{t\mathrm{\Pi }dtG_+}W)`$.
Now it is a starightforward calculation to check that the forms $`\stackrel{~}{C}_{g,n}^{ind}`$ (or $`C_{g,n}^{ind}`$) are $`(d+Q)`$-closed and satisfy the factorization property when restricted to the strata $`\gamma \times \{+\mathrm{}\}`$.
This construction is not smooth and is defined not on $`\overline{๐ฆ}_{g,n}`$, but on its extension. Nevertheless, one can easily turn this into a clear mathematical theory. We sketch the required construction in the next subsection.
#### 2.4.5. Moduli spaces with cuffs
Instead of Zwiebach invariants on the spaces $`\overline{๐ฎ}_{g,n}`$ we can consider Zwiebach invariants on the moduli spaces with cuffs. That is, at each boundary stratum of $`\overline{๐ฎ}_{g,n}`$ of codimension $`1`$ we glue the cylinder equal to this stratum multiplied by $`[0,+\mathrm{}]`$. Then we consider the set of forms satisfying the same axioms as above, but we require the properties of horizontality and factorization on the โ$`+\mathrm{}`$โ ends of the glued cylinders.
Thus we obtain a slight generalization of the notion of Zwiebach invariants. If we have a system of Zwiebach invariants on $`\overline{๐ฎ}_{g,n}`$, then we can lift it to cuffs. We just take the pull-backs of these forms under the mapping that keeps the moduli spaces and projects all cylinders to their โ$`0`$โ ends.
Then, when we consider the induced Zwiebach invariants, we glue new cylinders to the โ$`\mathrm{}`$โ ends of the cuffs. Thus, at each boundary stratum we have two consequently glued cylinders. So, we choose a certain mapping, which identifies two glued cylinders with one cylinder. Then the theory of induced Zwiebach invariants is again the theory of Zwiebach invariants on the moduli spaces with cuffs.
#### 2.4.6. Hodge case
In the Hodge case, we assume that $`QH^{}=G_{}H^{}=0`$. Then the induced Zwiebach invariants determine Gromov-Witten invariants obtained by integrals over the fundamental cycles. What we get is a sum over graphs with vertices marked by the initial Zwiebach invariants (or rather their integrals over the fundamental cycles), internal edges correspond to the contraction of outputs with the scalar product $`(,G_{}G_+)`$, and tails are marked by the elements of $`H^{}`$.
In this paper, we study the case where the unique nonvanishing integral of the initial Zwiebach invariants over the fundamental cycles exists for $`g=0`$, $`n=3`$. There are some obstruction for the existence of such initial Zwiebach invariants. We study them in the next subsection
#### 2.4.7. Obstructions
First, we choose $`\stackrel{~}{C}_{0,3}`$. It is a $`H^3`$-valued constant, so it determines a commutative multiplication on $`H`$. Since $`\stackrel{~}{C}_{0,3}`$ is $`Q`$-closed, we have the Leibnitz rule: $`Q(ab)=Q(a)b+aQ(b)`$.
Now we try to choose $`\stackrel{~}{C}_{0,4}`$. From the factorization property, it follows that $`\stackrel{~}{C}_{0,4}`$ is a sum of $`0`$-form and $`1`$-form. It is $`(d+Q)`$-exact. So, comparing values of the $`0`$-form at two different boundary cycles of $`\overline{๐ฆ}_{0,4}`$, we obtain that our multiplication is homotopy associative, that is, $`(ab)ca(bc)Q(H)`$.
Another relation comes from an attempt to glue the $`1`$-forms arising on the boundary of $`\overline{๐ฆ}_{0,4}`$ due to the factorization property. Consider the total space $`\overline{๐ฎ}_{0,4}`$. There are $`7`$ distinguished $`1`$-cycles, determined by the action of $`U(1)`$ at marked point and at double points. If we take an $`abcd`$-valued component of $`C_{g,n}`$, then from the factorization property, if follows that the integrals over these cycles are equal to $`(G_{}(a)b,cd),(G_{}(b)a,cd),(G_{}(c)d,ab),(G_{}(d)c,ab)`$ and $`(G_{}(ab),cd),(G_{}(ac),bd),(G_{}(ad),bc)`$ (the cycles corresponding to marked points are taken in the fiber over one of the boundary points). A path along each cycle can be obtained as a Dehn twist along the corresponding cycle on a surface with $`4`$ marked points. The relation among these Dehn twists imply that there is a $`2`$-dimensional surface in $`\overline{๐ฎ}_{0,4}`$, whose boundary is the sum of these seven $`1`$-cycles, and this gives us the $`7`$-term relation up to homotopy:
(22)
$$\begin{array}{c}G_{}(abc)+G_{}(a)bc+G_{}(b)ac+G_{}(c)ab\hfill \\ \hfill G_{}(ab)cG_{}(ac)bG_{}(bc)aQ(H)\end{array}$$
Now we try to choose $`\stackrel{~}{C}_{1,1}`$. The Dehn twists along the cycles on a genus $`1`$ surface with marked point also give us a new relation. There are three cycles, $`x`$ and $`y`$ are the basis in the first homology group of a torus, and $`z`$ is the cycle around the marked point. If $`D_x`$, $`D_y`$ and $`D_z`$ are the corresponding Dehn twists, then $`[D_x]=[D_y^1]`$ in the homology of $`\overline{๐ฆ}_{1,1}`$, and $`(D_yD_x^1D_y)^4=D_z`$ in the mapping class group . Therefore, we obtain that the kernel of the linear function
(23)
$$a(12str(G_{}a)str((G_{}a)))$$
contains the kernel of $`Q`$. Here $`str`$ denotes the supertrace, and $`a`$ (resp., $`(G_{}(a))`$) is the operator of multiplication by $`a`$ (resp., $`G_{}(a)`$).
From it follows that no other relations can come from the relations among Dehn twists. But of course there can be other obstructions of different geometric origin. We are grateful to E. Getzler for the explanation of the geometric origin of the $`7`$-term relation and $`1/12`$-axiom.
### 2.5. dGBV algebras
The simplest solutions to the relations presented in the previous subsection are known as differential Gerstenhaber-Batalin-Vilkovisky (dGBV) algebras, see . They have naturally appeared in the paper of Barannikov and Kontsevich as an axiomatic description of the properties of polyvector fields on Calabi-Yau. We refer to Pre-Introduction and to Appendix of this paper for the dicussion of additional benefits from the ideas hidden in .
We have seen above that it is very natural to obtain relations coming from the geometry of the moduli space of curves in calculations with graph constructions in dGBV algebras. In the rest of the paper we give a formal algebraic proof of WDVV and Getzler relations for the potential corresponding to simplest version of Zwiebach invarinats.
### 2.6. Acknowledgements
We are grateful to E. Getzler and M. Kontsevich for the fruitful discussions and to the referee, who has encouraged us to add the Pre-introduction and Appendix. Also, A. L. is grateful to A. Gerasimov for the explanation of the role of Hodge theory in string theory.
## 3. Construction
### 3.1. Hodge dGBV algebra
A Hodge differential Gerstenhaber-Batalin-Vilkovisky algebra is a supercommutative associative $``$-algebra $`H`$ with two odd linear operators
(24)
$$Q,G_{}:HH.$$
This operators must satisfy the system of axioms:
1. $`Q^2=G_{}^2=QG_{}+G_{}Q=0`$;
2. $`H=H_0H_4`$, where $`QH_0=G_{}H_0=0`$ and $`H_4`$ is represented as a direct sum of subspaces of dimension $`4`$ generated by $`e_\alpha ,Qe_\alpha ,G_{}e_\alpha ,QG_{}e_\alpha `$ for some vectors $`e_\alpha H_4`$, i. e.
(25)
$$H_4=\underset{\alpha }{}e_\alpha ,Qe_\alpha ,G_{}e_\alpha ,QG_{}e_\alpha ;$$
This axiom is called the axiom of Hodge decomposition. The ordinary dGBV-algebra is the structure that we have without axiom (2).
3. $`Q`$ is a derivation:
(26)
$$Q(ab)=Q(a)b+(1)^{\stackrel{~}{a}}aQ(b);$$
Here and below, we denote by $`\stackrel{~}{a}`$ the parity of $`aH`$.
4. $`G_{}`$ is an operator of the second order:
(27)
$$\begin{array}{c}G_{}(abc)=G_{}(ab)c+(1)^{\stackrel{~}{b}(\stackrel{~}{a}+1)}bG_{}(ac)+(1)^{\stackrel{~}{a}}aG_{}(bc)\hfill \\ \hfill G_{}(a)bc(1)^{\stackrel{~}{a}}aG_{}(b)c(1)^{\stackrel{~}{a}+\stackrel{~}{b}}abG_{}(c).\end{array}$$
Equation (27) is called the $`7`$-term relation.
### 3.2. Some notations
We define an operator $`G_+:HH`$. We set $`G_+H_0=0`$. On each subspace $`e_\alpha ,Qe_\alpha ,G_{}e_\alpha ,QG_{}e_\alpha `$, we define $`G_+`$ as $`G_+e_\alpha =G_+G_{}e_\alpha =0`$, $`G_+Qe_\alpha =e_\alpha `$, and $`G_+QG_{}e_\alpha =G_{}e_\alpha `$.
Clearly, $`G_+`$ is an odd operator, $`G_{}G_++G_+G_{}=0`$, and $`\mathrm{\Pi }_4=QG_++G_+Q`$ is the projection to $`H_4`$ along $`H_0`$. Denote by $`\mathrm{\Pi }_0`$ the projection to $`H_0`$ along $`H_4`$.
Thus $`G_+`$ is the homotopy operator corresponding to the contraction of $`H`$ to $`H_0`$. Note that we assume that this homotopy commutes with $`G_{}`$.
### 3.3. Integral
Let $`H`$ be a Hodge dGBV algebra. An integral on $`H`$ is an even linear function $`:H`$ such that
(28) $`{\displaystyle Q(a)b}`$ $`=(1)^{\stackrel{~}{a}+1}{\displaystyle aQ(b)},`$
(29) $`{\displaystyle G_{}(a)b}`$ $`=(1)^{\stackrel{~}{a}}{\displaystyle aG_{}(b)},`$
and
(30) $`{\displaystyle G_+(a)b}`$ $`=(1)^{\stackrel{~}{a}}{\displaystyle aG_+(b)}.`$
These properties imply that $`G_{}G_+(a)b=aG_{}G_+(b)`$, $`\mathrm{\Pi }_4(a)b=a\mathrm{\Pi }_4(b)`$, and $`\mathrm{\Pi }_0(a)b=a\mathrm{\Pi }_0(b)`$.
We define a scalar product on $`H`$:
(31)
$$(a,b)=ab.$$
We assume that this scalar product is non-degenerate.
We call the full structure that we have here (a Hodge dGBV algebra and an integral determining a non-degenerate scalar product on $`H`$) a *cyclic Hodge dGBV algebra*, or *cH algebra* for short. Further properties of this structure can be found in .
We would like to make two remarks on the scalar product (31):
1. Obviously, $`H_0`$ is orthogonal to $`H_4`$.
2. Using the non-degenerate scalar product (31), we may turn an operator $`A:HH`$ into the bivector (by bivector we call, for short, any element in $`H^2`$). Below we denote this bivector by $`[A]`$.
### 3.4. Variables
Let $`H_0`$ be a finite dimensional space. Let $`e_1,\mathrm{},e_n`$ be its basis. Denote by $`T_1,\mathrm{},T_n`$ some independent variables. We take the parity of $`T_i`$ equal to the parity of $`e_i`$.
### 3.5. Construction of potential
We construct a formal power series $`F=F_0+F_1+F_2+\mathrm{}`$ in variables $`T_1,\mathrm{},T_n`$.
We consider all trivalent graphs. This means that we consider graphs with vertices of index $`3`$ only and with possible half-edges (leaves). We mark all leaves by elements from the set $`L=\{e_1T_1,\mathrm{},e_kT_k\}`$.
We associate to each internal edge of a graph the bivector $`[G_{}G_+]`$. In our pictures, we denote this by thick black points on the edges. Each internal vertex (of index $`3`$) corresponds to the $`3`$-form $`m(a,b,c)=abc`$.
Now each graph gives us a monomial in $`T_1,\mathrm{},T_n`$ as follows. At each vertex we have three incoming edges. They give three inputs for the corresponding $`3`$-form $`m`$. Such input is either a โhalfโ of $`[G_{}G_+]`$<sup>3</sup><sup>3</sup>3We note that from Section 3.3, it follows that this bivector is symmetric. or an element of $`L`$. We take the product of values of $`3`$-forms $`m`$ on their inputs at all vertices of a graph. This is the monomial that we associate to the graph.
We take each graph with the combinatorial coefficient that is equal to the inverse order of its group of automorphisms.
Denote by $`J:HH`$ the operator $`J:h(1)^{\stackrel{~}{h}}h`$.<sup>4</sup><sup>4</sup>4In physics, this operator is known as the fermionic parity operator and is usually denoted by $`(1)^F`$. If we consider a graph with $`g`$ loops, then at $`g`$ edges we put the bivector $`[JG_{}G_+]`$ instead of $`[G_{}G_+]`$. These $`g`$ edges can be arbitrary ones, but with the only restriction: if we cut the graph at these edges, then we get a tree.
Thus we obtain a Feynman diagram expansion of the integral discussed in the Appendix.
### 3.6. Examples
We give some examples. Let $`a,b,c`$ be different elements of $`L`$. Consider the graph
(32)
$$\begin{array}{c}\text{}\end{array}.$$
The order of its group of automorphisms is equal to $`2`$. So, it gives the monomial
(33)
$$\begin{array}{c}\frac{1}{2}[G_{}G_+][G_{}G_+],(ab)(c)(ab)\hfill \\ \hfill =\frac{1}{2}abG_{}G_+\left(cG_{}G_+\left(ab\right)\right).\end{array}$$
Another example:
(34)
$$\begin{array}{c}\text{}\end{array}.$$
The order of its group of automorphisms of this graph is also equal to $`2`$. We have the monomial
(35)
$$\frac{1}{2}[JG_{}G_+],(a)=\frac{1}{2}str(G_{}G_+a).$$
Here $`str`$ is the supertrace. We recall that the supertrace of an operator $`A`$ is defined as $`str(A)=tr(JA)`$. Equation (35) means that this monomial is equal to the supertrace of the operator $`G_{}G_+a:HH`$, $`hG_{}G_+(ah)`$.
### 3.7. Potential
We denote by $`F`$ the formal sum of such monomials over all possible trivalent graphs with leaves marked by elements of $`L=\{e_1T_1,\mathrm{},e_kT_k\}`$. Of course, we identify isomorphic graphs.
$`F`$ is naturally represented as $`F_0+F_1+F_2+\mathrm{}`$, where $`F_i`$ is the sum over graphs with $`i`$ loops. We shall now draw the first few terms of $`F_0`$ and $`F_1`$. For brevity, we denote by $`E`$ the sum $`e_1T_1+\mathrm{}+e_kT_k`$.
(39) $`F_0`$ $`={\displaystyle \frac{1}{6}}\begin{array}{c}\text{}\end{array}+{\displaystyle \frac{1}{8}}\begin{array}{c}\text{}\end{array}+{\displaystyle \frac{1}{8}}\begin{array}{c}\text{}\end{array}+\mathrm{}`$
(43) $`F_1`$ $`={\displaystyle \frac{1}{2}}\begin{array}{c}\text{}\end{array}+{\displaystyle \frac{1}{4}}\begin{array}{c}\text{}\end{array}+{\displaystyle \frac{1}{4}}\begin{array}{c}\text{}\end{array}+\mathrm{}`$
## 4. WDVV equation
We consider the moduli space $`\overline{}_{0,4}`$. The cohomology classes of any two points of $`\overline{}_{0,4}`$ coincide. This gives a differential equation for the Gromov-Witten potential in genus zero. We check this differential equation in our construction.
### 4.1. Boundary points
We denote the classes of boundary points of $`\overline{}_{0,4}`$ by $`\mathrm{\Delta }_{12|34}`$, $`\mathrm{\Delta }_{13|24}`$, and $`\mathrm{\Delta }_{14|23}`$:
(44)
$$\begin{array}{ccc}\begin{array}{c}\text{}\end{array}& \begin{array}{c}\text{}\end{array}& \begin{array}{c}\text{}\end{array}\\ \mathrm{\Delta }_{12|34}& \mathrm{\Delta }_{13|24}& \mathrm{\Delta }_{14|23}\end{array}$$
We explain these pictures by the following example. The first picture denotes the moduli point of $`\overline{}_{0,4}`$ represented by a two-component curve such that the marked points $`1`$ and $`2`$ lie on one component and the marked points $`3`$ and $`4`$ lie on the other component.
We have $`\mathrm{\Delta }_{12|34}=\mathrm{\Delta }_{13|24}=\mathrm{\Delta }_{14|23}`$ in homology of $`\overline{}_{0,4}`$.
### 4.2. Differential equations
This relation gives us some differential equations. We suppose that $`F_0`$ is a formal power series in variables $`T_1,\mathrm{},T_n`$, and $`\eta _{ij}`$ is a metric on the space generated by $`T_1,\mathrm{},T_n`$. If all variables are even, we have:
(45)
$$\begin{array}{c}\frac{^3F_0}{T_1T_2T_k}\eta _{kl}\frac{^3F_0}{T_lT_3T_4}=\frac{^3F_0}{T_1T_3T_k}\eta _{kl}\frac{^3F_0}{T_lT_2T_4}\hfill \\ \hfill =\frac{^3F_0}{T_1T_4T_k}\eta _{kl}\frac{^3F_0}{T_lT_2T_3}.\end{array}$$
We have here three equations; each of them is called the Witten-Dijkgraaf-Verlinde-Verlinde (WDVV) equation.
### 4.3. Theorem
In our case, $`F_0`$ is the sum over trees. The metric $`\eta _{ij}`$ is given by the scalar product on $`H_0`$, $`\eta _{ij}=(e_i,e_j)`$.
###### Theorem 1.
$`F_0`$, $`\eta _{ij}`$ satisfy the WDVV equation.
We explain the simplest case of this theorem. Denote by $`[\mathrm{\Pi }_0]`$ the $`2`$-form corresponding to the operator $`\mathrm{\Pi }_0`$. We can put this bivector on an internal edge of a graph. We denote this by a thick white point on the edge. In the simplest case, the theorem states that the $`4`$-form
(46)
$$(t,u,v,w)\begin{array}{c}\text{}\end{array}$$
restricted to $`H_0`$ is symmetric. In other words, for any $`t,u,v,wH_0`$
(47)
$$tu\mathrm{\Pi }_0(vw)=tv\mathrm{\Pi }_0(uw)=tw\mathrm{\Pi }_0(uv).$$
We prove Theorem 1 in Section 9.4. The simplest case of Theorem 1 (given by Equation (47)) is discussed in detail in Section 7.
## 5. Getzler relation
Getzler elliptic relation is a linear relation among some natural complex codimension $`2`$ strata in the cohomology ring of the moduli space $`\overline{}_{1,4}`$. It gives a differential equation for Gromov-Witten potentials in genera zero and one. We prove that our construction satisfies this differential equation.
### 5.1. Cycles in $`\overline{}_{1,4}`$
We list the codimension two cycles entering Getzler relation.
(48)
$$\begin{array}{c}\begin{array}{cccc}\begin{array}{c}\text{}\end{array}& \begin{array}{c}\text{}\end{array}& \begin{array}{c}\text{}\end{array}& \begin{array}{c}\text{}\end{array}\\ \mathrm{\Delta }_{2,2}& \mathrm{\Delta }_{2,3}& \mathrm{\Delta }_{2,4}& \mathrm{\Delta }_{3,4}\end{array}\\ \begin{array}{ccc}\begin{array}{c}\text{}\end{array}& \begin{array}{c}\text{}\end{array}& \begin{array}{c}\text{}\end{array}\\ \mathrm{\Delta }_{0,3}& \mathrm{\Delta }_{0,4}& \mathrm{\Delta }_b\end{array}\end{array}$$
We use here the notations from . A line marked by $`1`$ corresponds to a genus one curve. An unmarked line corresponds to a genus zero curve. Notches correspond to the marked points.<sup>5</sup><sup>5</sup>5Note that here we use pictures with absolutely different meaning then in the rest of the paper. For instance, in all other pictures we put notches just to set operators on graphs.
For example, the generic point of the stratum $`\mathrm{\Delta }_{2,2}`$ is represented by a curve of genus one. It has no marked points, but it has two attached genus zero curves with two marked points on each of them.
Each picture means that we label marked points by the numbers $`\{1,2,3,4\}`$ in all possible ways. For example, there are $`3`$ variants for $`\mathrm{\Delta }_{2,2}`$ and $`12`$ variants for $`\mathrm{\Delta }_{2,3}`$.
### 5.2. Relation
Getzler elliptic relation:
(49)
$$12\mathrm{\Delta }_{2,2}4\mathrm{\Delta }_{2,3}2\mathrm{\Delta }_{2,4}+6\mathrm{\Delta }_{3,4}+\mathrm{\Delta }_{0,3}+\mathrm{\Delta }_{0,4}2\mathrm{\Delta }_b=0.$$
We rewrite this relation as a differential equation for the formal power series $`F_0`$ and $`F_1`$. If all variables are even, we have:
(50) $`\mathrm{\Delta }_{2,2}`$ $`{\displaystyle \frac{^3F_0}{T_1T_2T_i}}\eta _{ij}{\displaystyle \frac{^2F_1}{T_jT_k}}\eta _{kl}{\displaystyle \frac{^3F_0}{T_lT_3T_4}}`$
$`+{\displaystyle \frac{^3F_0}{T_1T_3T_i}}\eta _{ij}{\displaystyle \frac{^2F_1}{T_jT_k}}\eta _{kl}{\displaystyle \frac{^3F_0}{T_lT_2T_4}}`$
$`+{\displaystyle \frac{^3F_0}{T_1T_4T_i}}\eta _{ij}{\displaystyle \frac{^2F_1}{T_jT_k}}\eta _{kl}{\displaystyle \frac{^3F_0}{T_lT_2T_3}},`$
(51) $`\mathrm{\Delta }_{2,3}`$ $`{\displaystyle \frac{^2F_1}{T_1T_i}}\eta _{ij}{\displaystyle \frac{^3F_0}{T_jT_2T_k}}\eta _{kl}{\displaystyle \frac{^3F_0}{T_lT_3T_4}}`$
$`+11termsobtainedbypermutationsof\{1,2,3,4\},`$
$`\mathrm{}`$
(52) $`\mathrm{\Delta }_b`$ $`{\displaystyle \frac{^4F_0}{T_1T_2T_iT_k}}\eta _{ij}\eta _{kl}{\displaystyle \frac{^4F_0}{T_3T_4T_jT_l}}`$
$`+2termsobtainedbypermutationsof\{1,2,3,4\}.`$
### 5.3. The $`1/12`$-axiom
In our construction, $`F_0`$ is the sum over trees, $`F_1`$ is the sum over graphs with one loop, and metric is just $`\eta _{ij}=e_ie_j`$.
###### Theorem 2.
$`F_0,F_1,\eta _{ij}`$ satisfy Getzler relation, if
(53)
$$\begin{array}{c}\text{}\end{array}=\frac{1}{12}\begin{array}{c}\text{}\end{array}.$$
We explain these pictures. On the left hand side, we mark the loop by $`G_{}`$. This means that we put on the loop the bivector $`[G_{}]`$. On the right hand side, we put $`G_{}`$ on the leaf and we have an empty loop. This means that we apply $`G_{}`$ to the input on the leaf and that we put the bivector $`[\mathrm{Id}]`$ on the loop.
In order to simplify the understanding and to explain our notations, we rewrite the $`1/12`$-axiom (53) in terms of tensors and in terms of supertraces. In terms of tensors, the $`1/12`$-axiom looks like
(54)
$$[JG_{}],h=\frac{1}{12}[J],G_{}(h).$$
In terms of supertraces, the $`1/12`$-axiom means
(55)
$$str(G_{}h)=\frac{1}{12}str(G_{}(h)).$$
So, this is just a rigid version of the axiom (23) obtained from the relation among Dehn twists in the fundamental group of $`\overline{๐ฆ}_{1,1}`$. In fact, one can include this additional axiom in the definition of cH-algebra, since it has the same status as, say, the $`7`$-term relation.
### 5.4. The simplest case
We describe the simplest case of Theorem 2. Let $`a,b,c,d`$ be elements of $`\{e_1T_1,\mathrm{},e_nT_n\}`$. At each picture below, we distribute $`a,b,c,d`$ among leaves in all possible ways (in other words, we put the sum $`a+b+c+d`$ at each leaf). Then we calculate $`\mathrm{\Delta }_{2,2},\mathrm{},\mathrm{\Delta }_b`$ according to our rules and check the relation (49).<sup>6</sup><sup>6</sup>6We would like to note that the computations hidden behind these words are rather hard.
(58) $`\mathrm{\Delta }_{2,2}`$ $`={\displaystyle \frac{1}{16}}\begin{array}{c}\text{}\end{array}+{\displaystyle \frac{1}{16}}\begin{array}{c}\text{}\end{array}`$
(61) $`\mathrm{\Delta }_{2,3}`$ $`={\displaystyle \frac{1}{4}}\begin{array}{c}\text{}\end{array}+{\displaystyle \frac{1}{4}}\begin{array}{c}\text{}\end{array}`$
(64) $`\mathrm{\Delta }_{2,4}`$ $`={\displaystyle \frac{1}{8}}\begin{array}{c}\text{}\end{array}+{\displaystyle \frac{1}{4}}\begin{array}{c}\text{}\end{array}`$
(66) $`\mathrm{\Delta }_{3,4}`$ $`={\displaystyle \frac{1}{4}}\begin{array}{c}\text{}\end{array}`$
(69) $`\mathrm{\Delta }_{0,3}`$ $`={\displaystyle \frac{1}{4}}\begin{array}{c}\text{}\end{array}+{\displaystyle \frac{1}{2}}\begin{array}{c}\text{}\end{array}`$
(72) $`\mathrm{\Delta }_{0,4}`$ $`={\displaystyle \frac{1}{16}}\begin{array}{c}\text{}\end{array}+{\displaystyle \frac{1}{4}}\begin{array}{c}\text{}\end{array}`$
(76) $`\mathrm{\Delta }_b`$ $`={\displaystyle \frac{1}{4}}\begin{array}{c}\text{}\end{array}+{\displaystyle \frac{1}{4}}\begin{array}{c}\text{}\end{array}+{\displaystyle \frac{1}{16}}\begin{array}{c}\text{}\end{array}`$
As usual, an internal vertex corresponds to the integral of all inputs, an edge with the thick black point corresponds to the bivector $`[G_{}G_+]`$, and an edge with the thick white point corresponds to the bivector $`[\mathrm{\Pi }_0]`$.
### 5.5. Proof
We explain the proof of Theorem 2 in Section 9. The simplest case of Theorem 2 is discussed in Section 8
## 6. Strategy of proofs
We prove our theorems in two steps. For each theorem, the first step is the simplest case of a theorem. For both our theorems, Theorem 1 and Theorem 2, it is the case of degree $`4`$ ($`4`$ marked points on a surface and $`4`$ leaves in a graph).
Studying Gromov-Witten invariants, it is enough to have a relation in $`\overline{}_{0,4}`$ (or $`\overline{}_{1,4}`$) to prove a differential equation in any degree. Indeed, a relation in $`\overline{}_{0,4}`$ ($`\overline{}_{1,4}`$) can be lift to any $`\overline{}_{0,n}`$ ($`\overline{}_{1,n}`$), $`n4`$ via the projection forgetting all but four marked points. It is not the case in our construction.
Nevertheless, we have a general technique that allows us to extend an argument proving the simplest case of any relation to the argument that proves the corresponding differential equation in any degree.
So, our proofs are organized in three sections. First, we prove the simplest case of Theorem 1; second, we prove the simplest case of Theorem 2; third, we explain how one can extend our arguments to have the full proofs.
## 7. The simplest case of Theorem 1
For the convenience of the reader, we explain the proof of the simplest case of Theorem 1 in terms of tensor and in terms of graphs simultaneously. This gives also a number of illustrations to the correspondence between the language of graphs and the language of tensors.
### 7.1. The simplest case
We formulate the simplest case of Theorem 1. Consider $`a,b,c,dL=\{e_1T_1,\mathrm{},e_kT_k\}`$. Theorem 1 states that
(77)
is symmetric under premutations of $`a,b,c,d`$.
We prove this. We have the operator $`\mathrm{\Pi }_0`$ on the internal edge. Since $`\mathrm{\Pi }_0=\mathrm{Id}QG_+G_+Q`$, we have:
(78)
$$\begin{array}{c}\text{}\end{array}=\begin{array}{c}\text{}\end{array}\begin{array}{c}\text{}\end{array}\begin{array}{c}\text{}\end{array}.$$
Here we use a new object in our graphs, an internal vertex of index $`4`$. A vertex of index $`k`$ corresponds in our formulas to the $`k`$-form
(79)
$$m_k(a_1,\mathrm{},a_k)=a_1\mathrm{}a_k.$$
As usual, the inputs of this form correspond to the incoming edges and leaves.
So, Equation (78) can be rewritten just as
(80)
$$ab\mathrm{\Pi }_0(cd)=abcdabQG_+(cd)abG_+Q(cd).$$
Since $`Q(xy)=Q(x)y+(1)^{\stackrel{~}{x}}xQ(y)`$ and $`Q(x)y=(1)^{\stackrel{~}{x}+1}xQ(y)`$, we can write in terms of graphs that
(81)
$$\begin{array}{c}\text{}\end{array}+\begin{array}{c}\text{}\end{array}+\begin{array}{c}\text{}\end{array}=0$$
(it is the case of even inputs on leaves).
Thus we have:
(85) $`=\begin{array}{c}\text{}\end{array}+\begin{array}{c}\text{}\end{array};`$
(89) $`=\begin{array}{c}\text{}\end{array}+\begin{array}{c}\text{}\end{array}.`$
One can also rewrite these equations as <sup>7</sup><sup>7</sup>7Starting from here and up to the end of the paper we put the signs in formulas with graphs without any additional explanation. All signs in our formulas agree with each other. The choice of the sign at each picture is determined by the choice of the underlying tensor formula. So, we always put signs in the most convenient way, and one can check that the corresponding underlying tensor formulas agree with each other.
(90) $`{\displaystyle abQG_+(cd)}`$ $`={\displaystyle G_+(cd)Q(a)b}+{\displaystyle G_+(cd)Q(b)a};`$
(91) $`{\displaystyle abG_+Q(cd)}`$ $`={\displaystyle G_+(ab)Q(c)d}+{\displaystyle G_+(ab)Q(d)c}.`$
Since $`Qa=Qb=Qc=Qd=0`$, we have
(92)
$$\begin{array}{c}\text{}\end{array}=\begin{array}{c}\text{}\end{array}=0$$
and therefore
(93)
$$\begin{array}{c}\text{}\end{array}=\begin{array}{c}\text{}\end{array}.$$
The last expression is obviously symmetric under permutations of $`a`$, $`b`$, $`c`$, $`d`$. The simplest case of Theorem 1 is proved.
### 7.2. The next to the simplest case
We proceed to the next to the simplest case of Theorem 1. We ought to do it since it is not clear from the previous calculations how the full system of axioms of dGBV algebra is used.
Take $`a,b,c,d,eL=\{e_1T_1,\mathrm{},e_kT_k\}`$. Theorem 1 states that
(94)
$$\begin{array}{c}\begin{array}{c}\text{}\end{array}+\begin{array}{c}\text{}\end{array}+\begin{array}{c}\text{}\end{array}+\hfill \\ \hfill \begin{array}{c}\text{}\end{array}+\begin{array}{c}\text{}\end{array}+\begin{array}{c}\text{}\end{array}\end{array}$$
is symmetric under premutations of $`a,b,c,d`$.
We study the first summand of this expression. We have:
(95)
$$\begin{array}{c}\begin{array}{c}\text{}\end{array}=\hfill \\ \hfill \begin{array}{c}\text{}\end{array}\begin{array}{c}\text{}\end{array}\begin{array}{c}\text{}\end{array}.\end{array}$$
Since $`Q(ab)=Q(a)b+aQ(b)=0`$, the middle term of this expression in equal to $`0`$. For the last term, we have:
(96) $`Q\left(cG_{}G_+(de)\right)`$ $`=Q(c)G_{}G_+(de)+cQG_{}G_+(de)`$
$`=cG_{}QG_+(de)cG_{}G_+Q(de)`$
$`=cG_{}(de).`$
In particular, we use here $`\mathrm{\Pi }_4=QG_++G_+Q`$, $`G_{}\mathrm{\Pi }_4=G_{}`$, $`Q(de)=0`$.
This allows us to rewrite the Equation (95) as
(97)
$$\begin{array}{c}\text{}\end{array}=\begin{array}{c}\text{}\end{array}+\begin{array}{c}\text{}\end{array}.$$
In the same way we can write down the similar formulas for the next two summands of the Expression (94):
(101) $`=\begin{array}{c}\text{}\end{array}+\begin{array}{c}\text{}\end{array}`$
(105) $`=\begin{array}{c}\text{}\end{array}+\begin{array}{c}\text{}\end{array}`$
For $`G_{}`$ we can use the $`7`$-term relation (27). Note that $`G_{}(c)=G_{}(d)=G_{}(e)=0`$. This yields:
(106)
$$G_{}(cde)=G_{}(cd)e+G_{}(ce)d+G_{}(de)c$$
and therefore
(107)
$$G_+\left(G_{}(cd)e\right)+G_+\left(G_{}(ce)d\right)+G_+\left(G_{}(de)c\right)=G_{}G_+(cde).$$
Using this, we see that the sum of the last summands of Equations (97), (101), and (105) is equal to
(108)
$$\begin{array}{c}\text{}\end{array}$$
Thus we have that the first line of Expression (94) is equal to
(109)
$$\begin{array}{c}\text{}\end{array}+\begin{array}{c}\text{}\end{array}+\begin{array}{c}\text{}\end{array}\begin{array}{c}\text{}\end{array}.$$
The same argument proves that the second line of Expression (94) is equal to
(110)
$$\begin{array}{c}\text{}\end{array}+\begin{array}{c}\text{}\end{array}+\begin{array}{c}\text{}\end{array}\begin{array}{c}\text{}\end{array}.$$
Hence, Expression (94) is equal to
(111)
$$\begin{array}{c}\text{}\end{array}+\begin{array}{c}\text{}\end{array}+\begin{array}{c}\text{}\end{array}+\begin{array}{c}\text{}\end{array}.$$
Obviously, it is symmetric under permutations of $`a`$, $`b`$, $`c`$, $`d`$. The next to the simplest case of Theorem 1 is proved.
## 8. The simplest case of Theorem 2
We prove the simplest case of Theorem 2 in two steps. First, we represent each cycle as a linear combination of graphs $`P_1,\mathrm{},P_9`$:
$`P_1`$ $`=\begin{array}{c}\text{}\end{array}`$ $`P_2`$ $`=\begin{array}{c}\text{}\end{array}`$ $`P_3`$ $`=\begin{array}{c}\text{}\end{array}`$
$`P_4`$ $`=\begin{array}{c}\text{}\end{array}`$ $`P_5`$ $`=\begin{array}{c}\text{}\end{array}`$ $`P_6`$ $`=\begin{array}{c}\text{}\end{array}`$
$`P_7`$ $`=\begin{array}{c}\text{}\end{array}`$ $`P_8`$ $`=\begin{array}{c}\text{}\end{array}`$ $`P_9`$ $`=\begin{array}{c}\text{}\end{array}`$
Then we substitute these expressions into Getzler relation (49) and get zero.
### 8.1. The cycle $`\mathrm{\Delta }_{2,4}`$.
We recall that
(112)
$$\mathrm{\Delta }_{2,4}=\frac{1}{8}\begin{array}{c}\text{}\end{array}+\frac{1}{4}\begin{array}{c}\text{}\end{array}.$$
Here we put on leaves the sum $`e=a+b+c+d`$ of arbitrary four elements $`a,b,c,dL=\{e_1T_1,\mathrm{},e_kT_k\}`$.
Since $`\mathrm{\Pi }_0=\mathrm{Id}QG_+G_+Q`$, we have
(113)
$$\begin{array}{c}\begin{array}{c}\text{}\end{array}=\begin{array}{c}\text{}\end{array}\hfill \\ \hfill \begin{array}{c}\text{}\end{array}\begin{array}{c}\text{}\end{array}.\end{array}$$
Using Equation (81), we move $`Q`$ to the neighbouring edges. The third summand of the right hand side of Equation (113) is equal to zero. Indeed, we move $`Q`$ to leaves, and use that $`Q(e)=0`$. We consider the second summand of the right hand side of Equation (113). There we move $`Q`$ to the edge marked by $`\mathrm{\Pi }_0`$ and to the edge marked by $`G_{}G_+`$. In the first case we get zero, since $`Q\mathrm{\Pi }_0=0`$. In the second case, $`Q`$ transforms $`G_{}G_+`$ into $`G_{}`$ and goes to leaves (we do the same with the third summand of the right hand side of Equation (95)). Finally, we have
(114)
$$\begin{array}{c}\text{}\end{array}=\begin{array}{c}\text{}\end{array}+\begin{array}{c}\text{}\end{array}.$$
The same argument shows that
(115)
$$\begin{array}{c}\text{}\end{array}=\begin{array}{c}\text{}\end{array}+\begin{array}{c}\text{}\end{array}.$$
Thus, we have
(116)
$$\begin{array}{c}\mathrm{\Delta }_{2,4}=\frac{1}{8}\begin{array}{c}\text{}\end{array}+\frac{1}{4}\begin{array}{c}\text{}\end{array}\hfill \\ \hfill \frac{1}{8}\begin{array}{c}\text{}\end{array}+\frac{1}{4}\begin{array}{c}\text{}\end{array}.\end{array}$$
We consider the last two terms of this expression. We can apply here the $`7`$-term relation (27). Since $`G_{}\mathrm{\Pi }_0=0`$ and $`G_{}e=0`$, it takes the form
(117)
$$\frac{1}{8}\begin{array}{c}\text{}\end{array}+\frac{1}{4}\begin{array}{c}\text{}\end{array}=\frac{1}{8}\begin{array}{c}\text{}\end{array}$$
($`G_{}`$ jumps to the edge with $`G_+`$ and we get there $`G_{}G_+`$; exactly the same argument is used to obtain Equation (108)). Thus, we have
(118)
$$\mathrm{\Delta }_{2,4}=\frac{1}{4}\begin{array}{c}\text{}\end{array}.$$
Now we start the same procedure with the next thick white point. We have
(119)
$$\begin{array}{c}\begin{array}{c}\text{}\end{array}=\begin{array}{c}\text{}\end{array}\hfill \\ \hfill +\begin{array}{c}\text{}\end{array}+\begin{array}{c}\text{}\end{array}\end{array}$$
Applying the $`1/12`$-axiom (53), we have
(120)
$$\begin{array}{c}\text{}\end{array}=\frac{1}{12}\begin{array}{c}\text{}\end{array}.$$
From the $`7`$-term relation (27), it follows that $`G_{}(e^4)=2eG_{}(e^3)`$. Applying this, we have
(121)
$$\begin{array}{c}\text{}\end{array}=\frac{1}{2}\begin{array}{c}\text{}\end{array}.$$
So, the final formula for the cycle $`\mathrm{\Delta }_{2,4}`$ is
(122)
$$\mathrm{\Delta }_{2,4}=\frac{1}{4}\begin{array}{c}\text{}\end{array}\frac{1}{8}\begin{array}{c}\text{}\end{array}\frac{1}{48}\begin{array}{c}\text{}\end{array}.$$
### 8.2. The other cycles
The same calculations with the other cycles express these cycles in terms of the graphs $`P_1,\mathrm{},P_9`$:
$`\mathrm{\Delta }_{2,2}=`$ $`{\displaystyle \frac{1}{16}}P_1+{\displaystyle \frac{1}{16}}P_4{\displaystyle \frac{1}{8}}P_3+{\displaystyle \frac{1}{192}}P_9`$
$`\mathrm{\Delta }_{2,3}=`$ $`{\displaystyle \frac{1}{4}}P_1+{\displaystyle \frac{1}{4}}P_5{\displaystyle \frac{1}{4}}P_2`$
$`\mathrm{\Delta }_{2,4}=`$ $`{\displaystyle \frac{1}{4}}P_2{\displaystyle \frac{1}{8}}P_1{\displaystyle \frac{1}{48}}P_7`$
$`\mathrm{\Delta }_{3,4}=`$ $`{\displaystyle \frac{1}{4}}P_3{\displaystyle \frac{1}{12}}P_2{\displaystyle \frac{1}{48}}P_6+{\displaystyle \frac{1}{144}}P_7`$
$`\mathrm{\Delta }_{0,3}=`$ $`{\displaystyle \frac{1}{4}}P_6{\displaystyle \frac{1}{4}}P_8{\displaystyle \frac{1}{12}}P_7`$
$`\mathrm{\Delta }_{0,4}=`$ $`{\displaystyle \frac{1}{16}}P_9+{\displaystyle \frac{1}{4}}P_8{\displaystyle \frac{1}{8}}P_6`$
$`\mathrm{\Delta }_b=`$ $`{\displaystyle \frac{3}{8}}P_4{\displaystyle \frac{1}{2}}P_5+{\displaystyle \frac{1}{16}}P_9`$
Substituting these expressions into Getzler relation (49), we see that the coefficient at each $`P_i`$ is equal to zero. This proves the simplest case of Theorem 2.
## 9. General case of both Theorems
In this section, we will do the following. In order to prove our theorems in general case, we must consider graphs with an arbitrary number of leaves in addition to the basic four leaves that we consider in the simplest case. The idea is to use the โself-repeatingโ structure of our graphs. It means that we replace each edge marked by thick black point by the sum over all trivalent trees with two special leaves playing the role of the ends of the edge. In the similar way, we replace each edge marked by thick white point by the sum over all trivalent trees with two special leaves playing the role of the ends of the edge and a special edge marked by a thick white point on the path connecting these two leaves (all other edges are marked by thick black points, of course). Also we replace each leaf by the sum over rooted trivalent trees with a special leaf that corresponds to the initial one.
At the level of tensors this means that we replace in the formulas (66)-(77) for the simplest cases the operators $`\mathrm{\Pi }_0`$, $`G_{}G_+`$ and vectors $`a,b,c,d`$ by certain operators $`O_0`$, $`O_c`$, and vectors $`O_la,O_lb,O_lc,O_ld`$. We define all these operators ($`O_0,O_c`$, and $`O_l`$) in Section 9.2. In order to give compact definitions of these operators, we introduce in Section 9.1 an auxiliary vector $`\gamma `$ that is responsible, in a sense, for the self-repeating structure of our graphs. All our new operators, $`O_0,O_c`$, and $`O_l`$, are formal power seria in the variables $`T_1,\mathrm{},T_k`$. The degree zero part of these operators gives the simplest cases of our theorems. The degree one part of these operators gives the next to the simplest cases of our theorems. We give an example for this in Section 9.3. In Section 9.4 we complete the proof of Theorem 1, and in Section 9.5 we complete the proof of Theorem 2.
### 9.1. Vector $`\gamma `$
In this section, we define a vector
(123)
$$\gamma H[[T_1,\mathrm{},T_n]]$$
and study its properties. We denote by $`E`$ the sum $`E=e_1T_1+\mathrm{}+e_nT_n`$. We denote by $`\gamma `$ the outcome at the root of the sum of all rooted trivalent tries with $`E`$ on leaves and $`G_{}G_+`$ on edges:
(124)
$$\begin{array}{c}\gamma =\begin{array}{c}\text{}\end{array}+\frac{1}{2}\begin{array}{c}\text{}\end{array}+\frac{1}{2}\begin{array}{c}\text{}\end{array}\hfill \\ \hfill +\frac{1}{8}\begin{array}{c}\text{}\end{array}+\frac{1}{2}\begin{array}{c}\text{}\end{array}+\mathrm{}\end{array}$$
###### Lemma 1.
Vector $`\gamma `$ satisfied two equations:
(125) $`G_{}(\gamma )=0;`$
(126) $`Q(\gamma )+{\displaystyle \frac{1}{2}}G_{}(\gamma ^2)=0.`$
In particular, our $`\gamma `$ is a specific solution to the Maurer-Cartan equation defined in \[1, Lemma 6.1\]
We prove Lemma 1. The first statement is obvious, since $`G_{}E=0`$ and $`G_{}^2=0`$. We prove the second statement. Since $`[Q,G_{}G_+]=G_{}`$ and $`QE=0`$, and using the self-repeating structure of our graphs, we have:
(127)
$$Q(\gamma )=\frac{1}{2}\underset{i=0}{\overset{\mathrm{}}{}}\underset{i}{\underset{}{}}G_{}G_+\left(\gamma G_{}G_+\left(\gamma \mathrm{}G_{}G_+\left(\gamma G_{}\left(\gamma ^2\right)\right)\right)\right).$$
From the $`7`$-term relation (27), it follows that $`3\gamma G_{}(\gamma ^2)=G_{}(\gamma ^3)+3\gamma ^2G_{}(\gamma )=G_{}(\gamma ^3)`$, since $`G_{}(\gamma )=0`$. Substituting this in (127), we get
(128)
$$\begin{array}{c}Q(\gamma )=\frac{1}{2}G_{}(\gamma ^2)\hfill \\ \hfill \frac{1}{6}\underset{i=1}{\overset{\mathrm{}}{}}\underset{i1}{\underset{}{}}G_{}G_+\left(\gamma G_{}G_+\left(\gamma \mathrm{}G_{}G_+\left(\gamma G_{}G_+G_{}\left(\gamma ^3\right)\right)\right)\right).\end{array}$$
Since $`G_{}G_+G_{}=0`$, we have $`Q(\gamma )=(1/2)G_{}(\gamma ^2)`$. Lemma 1 is proved.
### 9.2. Some additional operators and vectors
In this section, we define some additional operators and vectors in terms of $`\gamma `$ and study their properties.
Define the operator $`\mathrm{\Gamma }`$,
(129)
$$\mathrm{\Gamma }(h)=G_{}G_+(\gamma h),$$
which obeys:
(130)
$$[Q,\mathrm{\Gamma }](h)=G_{}(\gamma h)G_{}G_+\left(\frac{\gamma ^2}{2}h\right).$$
Define the operator $`O_l`$ as:
(131)
$$O_l=\underset{i=0}{\overset{\mathrm{}}{}}\underset{i}{\underset{}{\mathrm{\Gamma }\mathrm{\Gamma }\mathrm{}\mathrm{\Gamma }}}.$$
Consider a vector $`aH_0C[[T_1,\mathrm{},T_k]]`$. Using Equation (130), Lemma 1, and the $`7`$-term relation (27), we have
(132)
$$QO_l(a)=G_{}\left(\gamma O_l(a)\right).$$
We will use the vector $`O_l(a)`$ instead of $`a`$ on leaves, and relation (132) instead of $`Qa=0`$. In terms of graphs the vector $`O_l(a)`$ can be represented as:
(133)
$$O_l(a)=\underset{i=0}{\overset{\mathrm{}}{}}\begin{array}{c}\text{}\end{array}$$
(the sum is taken over the number of fragments in graphs).
Define the operator $`O_c`$,
(134)
$$O_c=O_lG_{}G_+.$$
Using Equation (130), Lemma 1, and the $`7`$-term relation (27), we have
(135)
$$[Q,O_c](h)=G_{}\left(\gamma O_c(h)\right)O_c\left(\gamma G_{}(h)\right)G_{}(h).$$
We will use the operator $`O_c`$ instead of $`G_{}G_+`$ on edges, and relation (135) instead of $`[Q,G_{}G_+]=G_{}`$. We draw the operator $`O_c`$ in terms of graphs as:
(136)
$$O_c=\underset{i=0}{\overset{\mathrm{}}{}}\begin{array}{c}\text{}\end{array}$$
(the sum is taken over the number of fragments in graphs).
Define the operator $`O_r`$ as:
(137)
$$O_r(h)=h+\gamma O_lG_{}G_+(h).$$
Now consider the operator $`O_0`$ defined by the formula
(138)
$$O_0=O_l\mathrm{\Pi }_0O_r.$$
By applying several times Equation (130), Lemma 1, and the $`7`$-term relation (27), we arrive at:
(139)
$$O_0=O_l+O_r\mathrm{Id}[Q,O_lG_+O_r]+O_lG_+O_r\gamma G_{}G_{}\gamma O_lG_+O_r$$
(here we denote by $`\gamma `$ the operator of multiplication by $`\gamma `$). We will use the operator $`O_0`$ instead of $`\mathrm{\Pi }_0`$ on edges, and relation (139) instead of $`\mathrm{\Pi }_0=\mathrm{Id}QG_+G_+Q`$. We draw the operator $`O_0`$ in terms of graphs:
(140)
$$O_0=\underset{i,j=0}{\overset{\mathrm{}}{}}\begin{array}{c}\text{}\end{array}$$
(the sum is taken over the number of fragments and in graphs).
### 9.3. Degree one case
We study the case of degree one for Theorem 1. If we replace the operator $`\mathrm{\Pi }_0`$ by $`O_0`$ and the vectors $`a,b,c,d`$ by $`O_la,O_lb,O_lc,O_ld`$, then we have the following picture:
(141)
$$\begin{array}{c}\text{}\end{array}.$$
The operators $`O_l`$ and $`O_0`$ are the formal power seria in $`T_1,\mathrm{},T_k`$. We write down the first two terms of the power series expansions of these operators:
(142) $`O_l(x)`$ $`=\mathrm{Id}(x)+{\displaystyle \underset{i=1}{\overset{k}{}}}G_{}G_+(e_iT_ix)+\mathrm{}`$
(143) $`O_c(x)`$ $`=\mathrm{\Pi }_0(x)+{\displaystyle \underset{i=1}{\overset{k}{}}}\left(G_{}G_+(e_iT_i\mathrm{\Pi }_0(x))+\mathrm{\Pi }_0(e_iT_iG_{}G_+(x))\right)\mathrm{}`$
Then we have the power series expansion of picture (141)
(146) $`\begin{array}{c}\text{}\end{array}=`$
(149) $`+\begin{array}{c}\text{}\end{array}+\begin{array}{c}\text{}\end{array}`$
(152) $`+\begin{array}{c}\text{}\end{array}+\begin{array}{c}\text{}\end{array}`$
(155) $`+\begin{array}{c}\text{}\end{array}+\begin{array}{c}\text{}\end{array}`$
$`+\mathrm{}`$
Thus we see, that the degree zero part of the power series expansion of (141) is the simplest case of Theorem 1 (see Section 7), and the degree one part of it is the next to the simplest case of Theorem 1 (see Section 7.2).
### 9.4. Proof of Theorem 1
First we reformulate Theorem 1 in terms of $`O_0`$ and $`O_l`$. We claim that for any $`a,b,c,dL=\{e_1T_1,\mathrm{},e_kT_k\}`$,
(156)
is symmetric under permutations of $`a,b,c,d`$.
Using Equations (139) and (132) we can prove this exactly by the same argument as we prove the simplest case of this Theorem. Indeed, first we can use Equation (139) (instead of the formula $`\mathrm{\Pi }_0=\mathrm{Id}QG_+G_+Q`$). Using Equation (81), we have
(159) $`\begin{array}{c}\text{}\end{array}=`$ $`\begin{array}{c}\text{}\end{array}`$
(162) $`+\begin{array}{c}\text{}\end{array}+\begin{array}{c}\text{}\end{array}`$
(164) $`+\begin{array}{c}\text{}\end{array}`$
(166) $`+\begin{array}{c}\text{}\end{array}`$
(168) $`+\begin{array}{c}\text{}\end{array}`$
(170) $`+\begin{array}{c}\text{}\end{array}`$
(172) $`+\begin{array}{c}\text{}\end{array}`$
(174) $`+\begin{array}{c}\text{}\end{array}`$
(abusing notations, we denote by $`\gamma `$ the operator of multiplication by $`\gamma `$).
Applying the $`7`$-term relation (27) to the summands (164), (166), and (168) and using $`G_{}(\gamma )=G_{}(O_lc)=G_{}(O_ld)=0`$, we get that the sum of these three summands is equal to
(175)
$$\begin{array}{c}\text{}\end{array}.$$
Note that $`O_rG_{}=G_{}`$. Hence, $`O_lG_+O_rG_{}\gamma =O_lG_+G_{}\gamma =O_lG_{}G_+\gamma `$. Note also that $`O_lG_{}G_+\gamma =O_l\mathrm{Id}`$. Hence, the sum of (164), (166), and (168) is equal to
(176)
$$\begin{array}{c}\text{}\end{array}+\begin{array}{c}\text{}\end{array}.$$
The same argument proves that the sum of (170), (172), and (174) is equal to
(177)
$$\begin{array}{c}\text{}\end{array}+\begin{array}{c}\text{}\end{array}.$$
Substituting these expressions in Equation (159), we have
(178)
$$\begin{array}{c}\text{}\end{array}=\begin{array}{c}\text{}\end{array}.$$
The right hand side here is obviously symmetric under permutations of $`a,b,c,d`$. This proves Theorem 1.
### 9.5. On Theorem 2
We do not give here the detailed calculation proving the general case of Theorem 2. We just explain how to do this. It is obvious that our argument works, and calculations with Theorem 1 completely explain us what to do.
In order to have the full statement of Theorem 2, we change the markings on edges and leaves in pictures of the cycles $`\mathrm{\Delta }_{2,2},\mathrm{},\mathrm{\Delta }_b`$. So, we change $`\mathrm{\Pi }_0`$ to $`O_0`$, we change $`G_{}G_+`$ to $`O_c`$, and we change $`e`$ on leaves to $`Q_le`$.
In order to prove Theorem 2, we express these new cycles $`\mathrm{\Delta }_{2,2},\mathrm{},\mathrm{\Delta }_b`$ in terms of graphs $`P_1,\mathrm{},P_9`$, where we also change $`G_{}G_+`$ to $`O_c`$ and $`e`$ to $`Q_le`$.
Our calculations are just the same (like in the case of Theorem 1). But instead of the relation $`\mathrm{\Pi }_0=\mathrm{Id}QG_+G_+Q`$ we use Equation (139), instead of $`[Q,G_{}G_+]=G_{}`$ we use Equation (135), and instead of $`Qe=0`$ we use Equation (132).
The expressions of cycles $`\mathrm{\Delta }_{2,2},\mathrm{},\mathrm{\Delta }_b`$ in terms of graphs $`P_1,\mathrm{},P_9`$ are just the same as in the simplest case. Moreover, the intermediate step (Equation (118) for $`\mathrm{\Delta }_{2,4}`$) in calculations with each cycle is just the same as in the simplest case, but we must also change $`\mathrm{\Pi }_0`$, $`G_{}G_+`$, and $`e`$ to $`O_0`$, $`O_c`$, and $`O_le`$ in the intermediate pictures.
Finally, this proves Theorem 2. We note that this argument works not only for Theorem 1 and Theorem 2. That is, if we have any PDE for our potential $`F`$, which is proved in its simplest case by the same argument as we have used for the simplest cases of Theorem 1 and Theorem 2 (to get out step by step of thick white points increasing the indices of vertices), then the argument described here immediately gives the full proof of this PDE. This corresponds in the theory of Gromov-Witten invariants to the lift of relations among strata in the moduli spaces of curves (for example, Getzler relation in $`\overline{}_{1,4}`$ gives us relations in $`\overline{}_{1,5}`$, $`\overline{}_{1,6}`$, and so on).
## Appendix A BCOV-action
In this appendix, we explain how one can reformulate the results of Barannikov and Kontsevich in terms of graphs just by studying the BCOV-action proposed in the Appendix of their paper .
### A.1. Sums over trees
Let $`V`$ be an arbitrary vector space. Our goal is to find a critical point and the critical value at this point of the following expression:
(179)
$$A(v)=K_1(v)+\frac{1}{2}K_2(v,v)+\frac{1}{6}K_3(v,v,v)\frac{1}{2}B_2(v,v)$$
Here $`K_1`$, $`K_2`$, and $`K_3`$ are certain symmetric $`1`$-, $`2`$-, and $`3`$-forms respectively, and $`B_2`$ is a nondegenerate scalar product. We denote by $`b_2`$ the inverse bivector of $`B_2`$. Our goal is to obtain a critical point of $`A(v)`$ and the critical value at this point as a formal power series in $`K_i`$.
We consider the sum of rooted trees without leaves. We suppose that there are vertices of degree $`1`$, $`2`$, and $`3`$, and the root is the vertex of degree $`1`$. At each vertex (except the root) of degree $`i`$ we put the $`i`$-form $`K_i`$. At each edge we put the bivector $`b_2`$. Then, substituting the bivectors into the form according to the graph, we get a vector at the root. We also weight each graph with the inversed order of its automorphism group.
We denote the vector represented in this way by $`v_{cr}`$ (we suppose that the sum over rooted trees converges).
###### Lemma 2.
$`v_{cr}`$ is a critical point of $`A(v)`$.
Now we consider the sum over trees without leaves and without a root. We suppose that there are vertices of degree $`1`$, $`2`$, and $`3`$, and the root is the vertex of degree $`1`$. At each vertex of degree $`i`$ we put the $`i`$-form $`K_i`$. At each edge we put the bivector $`b_2`$. Substituting the bivectors into the form according to the graph, we get a number. As usual, we weight each graph with the inversed order of its automorphism group.
We denote the number obtained in this way by $`A_{cr}`$ (here we also suppose that the sum over trees converges).
###### Lemma 3.
$`A_{cr}=A(v_{cr})`$.
Both lemmas can be proved directly, by a simple linear algebra argument.
### A.2. BCOV-action
We consider a cH-algebra $`H`$. Barannikov and Kontsevich propose to study the action:
(180)
$$A(v)=\frac{1}{6}(E+G_{}v)^3\frac{1}{2}QvG_{}v.$$
We recall that $`E=e_1T_1+\mathrm{}+e_nT_n`$, $`n=dimH_0`$.
This is an immediate generalization of the Kodaira-Spenser theory of Bershadsky, Cecotti, Ooguri, and Vafa. However, the $`1/12`$-axiom is missing in and in all subsequent papers .
###### Proposition 1.
If $`v_{cr}`$ is the critical point of $`A(v)`$, then $`\gamma =E+G_{}(v_{cr})`$ is the $`G_{}`$-closed solution of the Maurer-Cartan equation (see Equations (125)-(126), Section 9.1).
###### Proposition 2.
The critical value $`F_0=A(v_{cr})`$ is the solution of the WDVV equation (see Equation (39), Section 3.7).
Barannikov and Kontsevich formulate and prove both propositions without using the representations of $`\gamma `$ and $`F_0`$ in terms of graphs. However, these representations exist and are naturally provided by the linear algebra formalism explained in the Section A.1.
Let us demonstrate this. The graph representation of $`F_0`$ is a direct corollary of the graph representation of $`\gamma `$. In order to obtain the graph representation of $`\gamma `$, we rewrite $`A(v)`$ as
(181)
$$\begin{array}{c}A(v)=\frac{1}{6}E^3+\frac{1}{2}E^2G_{}(v)\hfill \\ \hfill +\frac{1}{2}EG_{}(v)^2+\frac{1}{6}G_{}(v)^3\frac{1}{2}QvG_{}v.\end{array}$$
We recall that $`H=H_0_\alpha e_\alpha ,Qe_\alpha ,G_{}e_\alpha ,QG_{}e_\alpha `$. In fact, the scalar product $`B_2(v,v)=QvG_{}v`$ is nondegenerate only on $`_\alpha e_\alpha `$. So there exists the bivector $`b_2`$ inversed to $`B_2`$. We note that if we apply $`G_{}`$ to the both components of $`b_2`$, then we obtain the bivector $`[G_{}G_+]`$.
Now we consider the sum over the rooted trees discussed in the Section A.1. We have: $`K_1(v)=(1/2)E^2G_{}(v)`$, $`K_2(v,v)=EG_{}(v)^2`$, and $`K_3(v,v,v)=G_{}(v)^3`$. We see that we can move $`G_{}`$ from vertices to edges. Then, if we consider the sum over rooted trees with one additional $`G_{}`$ at the root, we obtain the following:
1. At edges we put the bivector $`[G_{}G_+]`$.
2. At vertices of degree $`3`$ we put the $`3`$-form $`(v_1,v_2,v_3)v_1v_2v_3`$.
3. At vertices of degree $`2`$ we put the $`2`$-form $`(v_1,v_2)v_1v_2E`$, i.e. we view it as the vertex of degree $`3`$ with one leaf marked by $`E`$.
4. At vertices of degree $`1`$ we put the $`1`$-form $`v_1v_1E^2/2`$, i.e. we view consider it as the vertex of degree $`3`$ with two leaves marked by $`E`$.
Thus we represented $`G_{}(v_{cr})`$ as a sum over the trivalent rooted trees with leaves. Moreover, $`E+G_{}(v_{cr})`$ is exactly the vector $`\gamma `$ studied in Section 9.1.
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# Contents
## 1 Introduction and motivation
In the presentโday analysis of various physical processes the complexity of the problem often makes it impossible to find a complete mathematical description of the experimentally measured or observed quantities. In order to proceed one usually simplifies the mathematical model which describes the bulk of the phenomena. The more delicate effects are accounted for through an expansion in terms of a small parameter which has to be defined such that one retains full quantitative control in the analysis. The art of a physicist consists in splitting the whole analysis into a main part and into small corrections to the main part for which one develops a perturbative expansion. One may quote one of the greatest physicists of the last century, L.D. Landau, who used to say that physics starts when a small parameter of the problem has been found. The usefulness and power of such an approach has been proven over and over for many years. An example that comes into mind are the first accurate calculations analyzing high precision data from celestial mechanics in terms of Newtonโs law of gravitation. When observational data on the motion of the planets became more accurate it was no longer sufficient to take into account only their gravitational interaction with the sun but it was necessary also to include the small corrections arising from the mutual gravitational interactions between the planets themselves. This was done through perturbation theory since the full equations for the planetary movements with mutual interaction fully taken into account are too difficult to treat analytically even within standard classical mechanics. A remarkable fact is that these corrections are very small but still have to be taken into account since the data are extremely precise. This is well illustrated by a historical example taken from astronomy which was the most precise scientific discipline in the past: the planets of the solar system close to Earth were directly observed while an eighth planet โ Neptune โ was first introduced in order to explain slight perturbations in the orbit of Uranus. The orbit of Neptune was theoretically calculated with high precision and the planet was directly observed later at the predicted place.
Let us now turn to high energy physics. The main interest of fundamental research in high energy physics is to find new physical phenomena and to place them into a consistent logical picture of the natural world which mankind inhabits. It is generally believed that the structure of fundamental interactions (except gravity) is qualitatively understood at the energies available at present accelerators. Basically there are two trends for the search of new physics. One is aiming at a direct discovery of new physical phenomena and fundamental constituents of matter (elementary particles) by moving to higher and higher energies. This requires ever more powerful accelerating facilities while the precision of the experimental detection and theoretical calculations can remain at a moderate level. The other is an indirect search for new physics at low energies (mainly through radiative corrections) which is based on very precise experimental data at moderate energies and a high accuracy of the theoretical calculations . A very recent example of the second category is the ongoing discussion about the possible appearance of deviations from the Standard Model (SM) of strong and electroweak interaction in new data on the muon anomalous magnetic moment . This two-fold road to scientific discovery is well known and has been pursued for ages as illustrated by the above example from celestial mechanics.
The muon anomalous magnetic moment has been computed in Quantum Electrodynamics (QED) โ the most precise physical theory of the present time. QED is a highly accurate theory of electromagnetic particle interactions where predictions for physical observables are calculated within perturbation theory in the form of a series in the expansion parameter of the theory, the fine structure constant $`\alpha `$ with a numerical value of $`\alpha ^1=137.036`$ at zero momentum. The fact that $`\alpha `$ is rather small makes the expansions well convergent. The muon anomalous magnetic moment has been computed through the expansion in the small ratio of lepton masses and the expansion in the coupling constant up to the order $`\alpha ^4`$. The most recent result reads
$$\alpha _\mu (\mathrm{QED})=A_1+A_2(m_\mu /m_e)+A_2(m_\mu /m_\tau )+A_3(m_\mu /m_e,m_\mu /m_\tau )$$
(1)
where for instance
$`A_1`$ $`=`$ $`0.5\left({\displaystyle \frac{\alpha }{\pi }}\right)0.328478965\mathrm{}\left({\displaystyle \frac{\alpha }{\pi }}\right)^2+1.181241456\mathrm{}\left({\displaystyle \frac{\alpha }{\pi }}\right)^3`$
$`1.5098(384)\left({\displaystyle \frac{\alpha }{\pi }}\right)^4+4.393(27)10^{12}.`$
The theoretical expressions for the coefficients of this series involve complicated multidimensional integrals the evaluation of which is the main difficulty in the calculation of a series expansion for the muon anomalous magnetic moment. A part of the coefficients in Eq. (1) have been calculated analytically whereas the numbers in the higher order terms have to be evaluated numerically .
With the advent of powerful computers a number of problems can now be approached through direct numerical calculations. The trajectories of satellites, for instance, are computed through direct calculations with computers. There are also important applications in Quantum Field Theory (QFT) especially, in the theory of strong interaction formulated on the lattice. However, at present the achieved accuracy is not yet sufficient leaving plenty of room for analytical computations. In addition, calculations in perturbation theory are often formulated in such a way that a numerical approach is not efficient or is unreliable due to the accumulation of rounding errors. The situation is quite analogous to the problem of long-time weather forecasts in meteorology where the direct numerical approaches are extremely unstable from the computational point of view. Rounding errors of the numerical evaluation of a large number of Feynman diagrams are such that the final accuracy is low because of huge cancellations between separate contributions of separate diagrams. Although an analytic evaluation is definitely preferable it cannot always be done due to technical difficulties in calculating the relevant Feynman integrals.
In Quantum Field Theory only a very limited number of problems have a full exact solution. Perturbation theory is a general tool for investigating the realistic physical situations. This is especially true in the theory of strong interactions and in the Standard Model in general. The classical examples are the analysis of $`e^+e^{}`$ annihilation into hadrons and the analysis of tau decays where perturbation theory is used in its most advanced form for the calculation up to very high orders of the strong coupling constant . In the latter case such an analysis has lead to the most accurate extraction of strong interaction parameters at low energies . The attempt to extrapolate the calculations to all orders for some limited subsets of diagrams or contributions are also quite useful . The coefficients of the expansions are given by complicated integrals which have a useful graphical representation in terms of Feynman diagrams. Graphical techniques are very efficient and convenient in the bookkeeping of the various loop contributions and are in wide use. In higher orders these graphs have many loops which are associated with the loop integrals. One speaks of multi-loop calculations for the evaluation of the coefficients of the perturbative expansion in QFT. Due to big advances in symbolic computing the most advanced symbolic programs can nowadays generate the terms of the perturbative expansion without actually drawing any graphs. Nevertheless, the terminology has survived. Nowadays the methods of QFT (and its terminology) is starting to penetrate into fields other than high energy physics such as polymer physics, hydro- and aerodynamics and weather forecasts which are much closer to real life. Therefore, the calculation of the complicated integrals associated with Feynman diagrams in a perturbative expansion is important in high energy physics as well as in many other areas of physics.
While the problem of calculating multi-loop integrals will not be solved in a general form in the foreseeable future, there is a continuing interest in the analysis of particular diagrams in the efforts to enlarge the number of known analytical results (as a review, see e.g. ). Especially promising in this regard is the analysis of subclasses of diagrams with simpler topological structures but with an arbitrary number of loops. The most promising in this respect is the simple topology of the so-called generalized sunrise diagrams (sometimes also called sunset, water melon, banana, or basketball diagrams). They will be referred to as the sunrise-type diagrams in this report. They can be computed for an arbitrary number of loops which makes them an ideal tool to obtain clues to the structure of more complicated multi-loop diagrams .
Recently there has been a renewal of interest in the calculation of diagrams of the sunrise-type topology for different numbers of loops involving different loop masses and/or external momenta of the diagram (see e.g. ). The classical (or genuine) two-loop sunrise diagram with different values of the internal masses has recently been studied in some detail using the traditional momentum space approach (see e.g. and references therein). In the present report we review the configuration space technique (also called $`x`$-space technique) for the calculation of sunrise-type diagrams. The configuration space technique is best suited for the topology of the diagram. It considerably reduces the complexity of the calculation and allows for new qualitative and quantitative results for this important particular class of Feynman integrals . Configuration space techniques can be used to verify known results obtained within other techniques both analytically and numerically and to investigate some general features of Feynman diagram calculation . We list features of the $`x`$-space technique in turn with the aim to show how efficient the method is.
The diagrams with sunrise-type topology form a subset of the general set of diagrams with a given number of loops. They appear in various specific physics applications:
* multi-loop calculations in general
* the analysis of static properties of baryons using sum rule techniques in QCD
* the properties of glueballs extracted from the perturbative analysis of two-point correlators
* the multi-particle phase space for multi-body decays of a single-particle phase space
* lattice QCD calculations
* mixing of neutral mesons
* Chiral perturbation theory (ChPT) and effective theories for Goldstone modes in higher orders of the momentum expansion
* analysis of exotic states in QFT: multi-quark states in QCD or pentaquarks using sum rule techniques
* general questions of QFT: sum rules in two-dimensional QED and properties of baryons in the large $`N_c`$ limit of QCD
* effective potentials for symmetry breaking
* finite temperature calculations
* applications in nuclear physics
* phase space integrals for particles in jets where the momentum along the direction of the jet is fixed
* applications in solid state physics
The principal aim of this report is to assemble all the necessary tools needed for the computation of multi-loop Feynman diagrams with sunrise-type topology. We discuss the whole known spectrum of different methods to analyze these integrals. Amomg these are concise analytical evaluation techniques, expansion in small parameters such as masses and momenta (or inverse masses and momenta), expansions in special kinematical regions as the threshold regime in the Minkowskian domain, integral representations for integrals including their discontinuity across physical cuts in terms of analytic functions of the external momentum, and, finally, efficient, fast, simple and stable numerical procedures. In the Euclidean domain the numerical procedures derived from our representation are efficient and reliable, i.e. stable against error accumulation. In order to acquaint the reader with the tools given given in the report we start with basic notations and relations.
We present a comprehensive report on configuration space techniques including many technical details such that an interested reader can use this report as a practical guide for practical calculations. We have included a great deal of mathematical material such that the report is self-contained to a large degree. In most cases it is not necessary to consult mathematical handbooks in order to understand the calculations. We also give a rather large sample of worked-out examples. The calculational methods used are rather well-suited for further development and can easily be tailored to the further specific needs of the potential user. Using the results of this paper one can create oneโs own software for an efficient evaluation of the many quantities of interest that can be extracted from the analysis of sunrise-type diagrams.
The report is organized as follows. In Sec. 2 we briefly summarize some general notions and their application to sunrise-type diagrams. In Sec. 2.1 we introduce the configuration space representation of sunrise-type diagrams and fix our notation. In Sec. 2.2 we specify what is meant by computing a sunrise-type diagram. In Sec. 2.3 we discuss the ultraviolet (UV) divergence structure of sunrise-type diagrams and present recipes to regularize the UV divergences either by subtraction or by dimensional regularization. In Sec. 2.4 we comment on the spectral density of sunrise-type diagrams and its connection to phase space. Sec. 3 is devoted to some explicit configuration space calculations involving both analytical and numerical methods the results of which are compared to previously known results where other calculational techniques have been used. One of the examples is the direct computation of the spectral density of sunrise-type diagrams without taking recourse to Fourier transforms. The results are important for multi-body phase space calculations. In Sec. 4 we discuss methods to find asymptotic expansions in different kinematical regimes of mass and/or momentum configurations in the Euclidean domain. Sec. 5 contains a generalization to non-standard propagators and non-scalar cases. In Sec. 6 we generalize the configuration space technique to other topologies. Sec. 7 contains our conclusions. Some of the lengthier formulas are relegated to the appendices, together with useful mathematical material about Bessel functions and Gegenbauer polynomials and a short treatise about cuts and discontinuities as they are occuring in the main text.
## 2 Basic notions and relations for sunrise-type diagrams
Sunrise-type diagrams are graphic representations of the $`n`$-loop two-point correlation functions in QFT with $`(n+1)`$ internal propagators connecting the initial and final vertex. The well-studied genuine sunrise diagram shown in Fig. 1(a) is the leading order perturbative correction to the lowest order propagator in $`\varphi ^4`$-theory, i.e. it is a two-point two-loop diagram with three internal lines. This diagram emerges as a correction to the Higgs boson propagation in the Standard Model. It also naturally appears in some effective theory approaches to critical phenomena and studies of phase transitions in QFT. The corresponding leading order perturbative correction in $`\varphi ^3`$-theory is a one-loop diagram which can be considered as a oversimplified case of the prior example.
The two-loop case is a standard starting point for the calculation of radiative corrections in QFT. It emerges in a huge variety of physical applications. Fortunately it can be analytically computed in any desired kinematical regime for any values of the relevant parameters. In this respect it does not present a real challenge as far as multi-loop calculations are concerned. We shall, nevertheless, discuss the genuine sunrise diagram in some detail in order to illustrate the efficiency of our configuration space methods. We shall compare our results with known exact analytical results obtained with other techniques. This provides a mutual check on the correctness of the results. A straightforward generalization of this topology to the multi-loop case is a correction to the free propagator in $`\varphi ^{n+2}`$-theory that contains $`n`$ loops and $`(n+1)`$ internal lines (see Fig. 1(b)).
A subclass of the general sunrise-type diagram shown in Fig. 1(b) is the case when the two external momenta vanish. This subclass is referred to as the subclass of vacuum bubble diagrams with sunrise-type topology which will be referred to as vacuum bubbles. They have a simpler structure than the true sunrise-type diagrams since the number of mass scales is reduced by one for the vacuum bubbles. We shall frequently return to the discussion of vacuum bubbles in the main text.
### 2.1 Definitions and notation
A $`n`$-loop sunrise-type diagram is a two-point correlation function with $`(n+1)`$ propagators connecting points $`x`$ and $`y`$ and as such is explicitly given by a product of simple propagators built from the basic two-point correlation function $`D(x,y,m)`$,
$$\mathrm{\Pi }(x,y)=\underset{i=1}{\overset{n+1}{}}D(x,y,m_i)$$
(2)
(and/or their derivatives if necessary). The quantity $`D(x,y,m)`$ may be the propagator of a free massive particle with mass $`m`$ or a more general two-point correlation function. The sunrise-type diagrams are the leading order contribution to a two-point correlation function of two local currents $`j_{n+1}(x)`$ of the form
$$j_{n+1}(x)=๐_{\mu _1}\varphi _1\mathrm{}๐_{\mu _{n+1}}\varphi _{n+1}$$
(3)
where the fields $`\varphi _i`$ have masses $`m_i`$ and where $`๐_\mu `$ is a derivative with multi-index $`\mu =\{\mu _1,\mathrm{},\mu _k\}`$ standing for $`๐_\mu =^k/x_{\mu _1}\mathrm{}x_{\mu _k}`$. The sunrise-type diagrams are contained in the leading order expression for the polarization function
$$\mathrm{\Pi }(x,y)=Tj_{n+1}(x)j_{n^{}+1}(y)$$
(4)
where the brackets mean quantum mechanical averaging over the ground state which is explicitly given by a product of propagators and/or their derivatives,
$$\mathrm{\Pi }(x,y)=๐_{\mu _1\nu _1}(x,y,m_1)\mathrm{}๐_{\mu _{n+1}\nu _{n+1}}(x,y,m_{n+1}).$$
(5)
As the standard case we will consider a translation invariant situation in which the propagator depends on the difference of the arguments only, $`D(x,y,m)=D(xy,m)`$. An exception to this important standard case would be the existence of an arbitrary external field. However, in the present report we will not further discuss this possibility. In the standard case one of the vertices can conveniently be placed at the origin, say $`y=0`$. The propagator is then a function of $`x`$ only, $`D(x,m)`$. The basic expression for the sunrise-type diagram in configuration space reads
$$\mathrm{\Pi }(x)=\underset{i=1}{\overset{n+1}{}}D(x,m_i).$$
(6)
$`D(x,m)`$ represents a free propagator of a massive particle with mass $`m`$ in $`D`$-dimensional (Euclidean) space-time. It is given by
$$D(x,m)=\frac{1}{(2\pi )^D}\frac{e^{i(px)}d^Dp}{p^2+m^2}=\frac{(mx)^\lambda K_\lambda (mx)}{(2\pi )^{\lambda +1}x^{2\lambda }}$$
(7)
where we have used $`D=2\lambda +2`$. $`K_\lambda (z)`$ is the modified Bessel function of the second kind (sometimes also known as McDonald function, see Appendix A) defined by Eq. (A8).
The massless propagator can be obtained from Eq. (7) by taking the limit $`m0`$ (or more precisely the limit $`mx0`$ at fixed $`x`$ in terms of the dimensionless quantity $`mx`$). It reads
$$D(x,0)=\frac{1}{(2\pi )^D}\frac{e^{i(px)}d^Dp}{p^2}=\frac{\mathrm{\Gamma }(\lambda )}{4\pi ^{\lambda +1}x^{2\lambda }}.$$
(8)
and where $`\mathrm{\Gamma }(\lambda )`$ is Eulerโs Gamma function.
In some physics applications one may have more general basic two-point functions (โmodified propagatorsโ) for massive particles. An example is the calculation of the three-point function necessary for the determination of a particle form factor at zero momentum transfer when the external momentum of the current vanishes. This momentum configuration reduces the diagram to a sunrise-type diagram with more complicated basic two-point lines (see e.g. ). This phenomenon frequently occurs also in calculations of various quantities in Chiral Perturbation Theory (ChPT) . Formally, these two-point functions cannot be considered as propagators of a physical particle but emerge as effective basic two-point functions. The simplest modification of the basic propagator which frequently appears in applications is the occurrence of higher powers of the standard propagator. In momentum space it has the form
$$\stackrel{~}{D}^{(\mu )}(p,m)=\frac{1}{(p^2+m^2)^{\mu +1}}.$$
(9)
For integer $`\mu `$ these cases appear if one e.g. considers mass derivatives of the propagator as they are needed in certain parameter expansions. For instance, one can consider
$$\stackrel{~}{D}^{}(p,m)=\frac{}{m^2}\frac{1}{(p^2+m^2)}$$
(10)
In the same way a derivative of the momentum itself will also increase the power of the propagator. For instance, the second order scalar derivative $`\text{ }\text{ }\text{ }=_\mu ^\mu `$ of the standard propagator will result in
$$_\mu ^\mu \frac{1}{p^2+m^2}=\frac{2}{(p^2+m^2)^2}+\frac{8p^2}{(p^2+m^2)^3}.$$
(11)
Higher powers of propagators occur in many applications. Therefore, any calculational technique should be well suited to treat such higher powers. It turns out that the configuration space method very naturally accommodates higher order derivatives.
In configuration space higher powers of the standard propagator with mass $`m`$ are explicitly given by (cf. Eq. (7))
$$\stackrel{~}{D}^{(\mu )}(x,m)=\frac{1}{(2\pi )^D}\frac{e^{i(px)}d^Dp}{(p^2+m^2)^{\mu +1}}=\frac{1}{(2\pi )^{\lambda +1}2^\mu \mathrm{\Gamma }(\mu +1)}\left(\frac{m}{x}\right)^{\lambda \mu }K_{\lambda \mu }(mx).$$
(12)
the additional power of denominator factors (or, equivalently, the number of derivatives in masses) is denoted in the corresponding Feynman diagram by dots placed on the line (though even the power can be non-integer). Note that the modified propagator has the same functional form as the usual propagator. The main difference is the change of the value of the index of the Bessel function $`K_\lambda (mx)K_{\lambda \mu }(mx)`$. This fact unifies and simplifies the analysis of sunrise-type diagrams that contain more complicated basic propagators (analytical expressions that correspond to one line of the diagrams) as they emerge in some physical applications or are used for the calculation of other diagrams through recurrence relations techniques .
### 2.2 Momentum versus configuration space representation: <br>Fourier transform in the evaluation of sunriseโtype integrals
The calculation of loop diagrams in perturbative Quantum Field Theory is associated with the evaluation of complicated multi-dimensional integrals, see e.g. . It is obvious that the expression for the sunrise-type diagram in configuration space given by Eq. (6) contains no integration at all even if it represents a multi-loop diagram. In this respect it is a kind of tree-level diagram that usually appears in the lowest orders of perturbative QFT. However, โto calculateโ a sunrise diagram means in most cases to find certain integrals of Eq. (6) over $`x`$ with some weight functions depending on the quantity in question. The simplest situation corresponds to a computation of the values of the vacuum bubbles, i.e. expressions for the diagrams without external momenta (or with vanishing external momenta, see e.g. ). In case of the sunrise-type topology the value of the vacuum bubble is given by integrals with a simple weight function,
$$\stackrel{~}{\mathrm{\Pi }}(0)=\mathrm{\Pi }(x)d^Dx=D(x,m_1)\mathrm{}D(x,m_{n+1})d^Dx.$$
(13)
This integral provides the simplest example of calculating a sunrise-type diagram. We like to emphasize that this calculation is almost always a must as it allows one to determine the counterterms related to UV divergences of sunrise-type diagrams that always occur in realistic physical models. Due to rotation invariance the integral is one-dimensional as the angular integration can be explicitly done. Thus
$`\stackrel{~}{\mathrm{\Pi }}(0)`$ $`=`$ $`{\displaystyle \mathrm{\Pi }(x)d^Dx}={\displaystyle D(x,m_1)\mathrm{}D(x,m_{n+1})d^Dx}`$ (14)
$`=`$ $`{\displaystyle \frac{2\pi ^{\lambda +1}}{\mathrm{\Gamma }(\lambda +1)}}{\displaystyle D(x,m_1)\mathrm{}D(x,m_{n+1})x^{2\lambda +1}๐x}.`$
The possibility to compute all vacuum bubbles in terms of one-dimensional integrals for any number of loops is a big advantage of the configuration space technique. This guarantees that the complete renormalization program can be explicitly accomplished in terms of one-dimensional integrals. It should be clear that this is a huge simplification.
The most familiar example for the evaluation of a sunrise-type diagram is the procedure of finding its value in momentum space, or put differently, the evaluation of its Fourier transform. It corresponds to non-bubble diagrams with two external lines as the integrals then requires the computation of the Fourier transform of $`\mathrm{\Pi }(x)`$,
$$\stackrel{~}{\mathrm{\Pi }}(p)=\mathrm{\Pi }(x)e^{i(px)}d^Dx.$$
(15)
The required integrals are basically scalar integrals. This makes the angular integration in Eq. (13) simple in $`D`$-dimensional space-time. The needed formula is the integration over the angular variables $`d^D\widehat{x}`$ which, in explicit form, reads
$$d^D\widehat{x}e^{i(px)}=2\pi ^{\lambda +1}\left(\frac{px}{2}\right)^\lambda J_\lambda (px)$$
(16)
where $`p=|p|`$, $`x=|x|`$, and $`J_\lambda (z)`$ is the Bessel function of the first kind. The integration $`d^D\widehat{x}`$ over the rotationally invariant measure $`d^D\widehat{x}`$ on the unit sphere in $`D`$-dimensional (Euclidean) space-time means integration over angles only. The generalization of Eq. (16) to more complicated integrands with an additional tensor structure $`x^{\mu _1}\mathrm{}x^{\mu _k}`$ is straightforward and leads to different orders $`J_{\lambda +l}(z)`$ of the Bessel function for different irreducible tensors (denoted by the number $`l`$) after angular averaging (integration over angles). The corresponding order of the Bessel function can be inferred from the expansion of the plane wave function $`e^{i(px)}`$ which appears as a weight function upon computing the Fourier transform in a series of Gegenbauer polynomials $`C_l^\lambda (w)`$ (see Appendix B). The Gegenbauer polynomials are orthogonal on a $`D`$-dimensional unit sphere, and the expansion of the plane wave in terms of Gegenbauer polynomials reads
$$e^{i(px)}=\mathrm{\Gamma }(\lambda )\left(\frac{px}{2}\right)^\lambda \underset{l=0}{\overset{\mathrm{}}{}}i^l(\lambda +l)J_{\lambda +l}(px)C_l^\lambda \left(\frac{(px)}{px}\right).$$
(17)
Besides the definition of the Gegenbauer polynomials by a characteristic polynomial as given in Eq. (B1), the expansion in Eq. (17) can also serve as a definition of these polynomials.
The representation in Eq. (17) allows one to single out an irreducible tensorial structure from the angular integration given by the Fourier integral in Eq. (13). Integration techniques involving Gegenbauer polynomials for the computation of multi-loop massless diagrams are described in detail in where many useful relations can be found (see also where the calculation in momentum space was considered). One arrives at a representation of the Fourier transform of a sunrise-type diagram that is given in terms of the one-dimensional integral
$`\stackrel{~}{\mathrm{\Pi }}(p)`$ $`=`$ $`2\pi ^{\lambda +1}{\displaystyle _0^{\mathrm{}}}\left({\displaystyle \frac{px}{2}}\right)^\lambda J_\lambda (px)\mathrm{\Pi }(x)x^{2\lambda +1}๐x`$ (18)
$`=`$ $`2\pi ^{\lambda +1}{\displaystyle _0^{\mathrm{}}}\left({\displaystyle \frac{px}{2}}\right)^\lambda J_\lambda (px)D(x,m_1)\mathrm{}D(x,m_{n+1})x^{2\lambda +1}๐x`$
which is simple to analyze. The representation given by Eq. (18) is quite universal.
### 2.3 Regularization and renormalization
While the configuration space expression for the sunrise-type diagram contains no integration it does not always represent an integrable function of $`x`$ in configuration space for general values of the space-time dimension $`D`$. The expression given by Eq. (6) can have non-integrable singularities at small $`x`$ for a sufficiently large number of propagators if the space-time dimension is $`D>2`$.
Upon multiplication of many propagators one may get a function which is singular at small $`x`$. Each propagator by itself can be considered as a distribution. However, distributions do not form an algebra and their multiplication is not well-defined . In this respect the polarization function $`\mathrm{\Pi }(x)`$ in Eq. (6) is not completely defined as a proper distribution. In attempting to integrate such a function over the whole $`x`$-space we encounter infinities in the form of ultraviolet (UV) divergences . Therefore, the computation of the Fourier transform of $`\mathrm{\Pi }(x)`$ requires regularization (for instance, dimensional regularization) and subtraction. Configuration space provides a nice environment for this.
Note that all these statements depend on the dimension of space-time. The vicinity of $`D=2`$ is special since the leading singularity of the basic propagator is only logarithmic and the multiplication of propagators does not lead to non-integrable singularities. Therefore, the UV properties of two-dimensional models that are often considered as simplified models in QFT or that emerge in real physical applications as a result of some given approximations . They are much softer than those of realistic four-dimensional models (see e.g. ).
In general, however, one has to consider an arbitrary space-time dimension $`D`$. The modern way of dealing with singularities in QFT in particular in multi-loop calculations is mostly based on dimensional regularization (as a review see ). Dimensional regularization is characterized by considering the space-time dimension $`D`$ to be an arbitrary parameter, for instance $`D=D_02\epsilon `$ where $`D_0`$ is an integer dimension.<sup>1</sup><sup>1</sup>1The usual choice of the integer space-time dimension $`D_0`$ for real physical applications is $`D_0=4`$, but other space-time dimensions $`D_0=2`$, $`D_0=3`$, $`D_0=5`$ or $`D_0=6`$ are being used in other applications. $`D_0=3`$ is natural for near threshold expansions or expansions in the Minkowskian domain in general. In case of non-integer space-time dimension $`D`$ the integration is well-defined, while the limit $`DD_0`$ is singular, consisting usually of poles in $`DD_0`$. These pole parts are extremely simple to handle in configuration space. Indeed, in configuration space the regularization is given by adding the $`\delta `$-function $`\delta (x)`$ and its derivatives to $`\mathrm{\Pi }(x)`$. As we will see in the following, this procedure corresponds to adding a polynomial in momentum space (local counterterms according to the Bogoliubov-Parasiuk-Hepp-Zimmermann (BPHZ) theorem ). The finite parts can be found numerically for any number of loops by doing a single one-dimensional integration. Therefore, the renormalization in $`x`$-space becomes
$$\mathrm{\Pi }^R(x)=\mathrm{\Pi }(x)+C_0\delta (x)+C_1\text{ }\text{ }\text{ }\delta (x)+\mathrm{}+C_r\text{ }\text{ }\text{ }^r\delta (x)$$
(19)
where $`\text{ }\text{ }\text{ }=_\mu ^\mu `$. The coefficients $`C_i`$ are functions of the regularization parameter $`\epsilon `$ which are singular in the limit $`\epsilon 0`$. The extraction of the coefficients $`C_i`$ is explained in the subsequent sections for some special cases. The number $`r`$ of counterterms is determined by dimensional arguments such as power counting and can be easily related to the number of propagators and the space-time dimension. The actual number of poles depends on the mass configuration.
### 2.4 The spectrum as discontinuity across the physical cut
From the physics point of view one important aspect of our analysis of sunrise-type diagrams is the construction of the spectral decomposition of the diagrams. For the two-point correlation function we determine the discontinuity across the physical cut in the complex plane of the squared momentum, $`p^2=m^2\pm i0`$ which is referred to as the spectral density
$$\rho (s)=\frac{1}{2\pi i}\mathrm{Disc}\stackrel{~}{\mathrm{\Pi }}(p)|_{p^2=s}$$
(cf. Appendix C). Note that the spectral density of a sunrise-type diagram is finite for any number of loops. It turns out that the configuration space technique allows one to compute the spectral density in a very efficient manner. The analytic structure of the correlator $`\mathrm{\Pi }(x)`$ (or the spectral density of the corresponding polarization operator) can be determined directly in configuration space without having to compute its Fourier transform first. The technique for the direct construction of the spectral density of sunrise-type diagrams introduced in is based on an integral transform in configuration space which in turn is given by the inversion of the relevant dispersion relation.
The dispersion representation (or spectral decomposition) of the polarization function in configuration space has the form
$$\mathrm{\Pi }(x)=_0^{\mathrm{}}\rho (s)D(x,\sqrt{s})๐s=_0^{\mathrm{}}\rho (m^2)D(x,m)๐m^2$$
(20)
where $`\sqrt{s}=m`$. This representation was used for sum rule applications in where the spectral density for the two-loop sunrise diagram was given in two-dimensional space-time . The representation in momentum space is more familiar and is referred to as the Kรคllรฉn-Lehmann representation of the two-point correlation function. In the Euclidean domain it is given by
$$\stackrel{~}{\mathrm{\Pi }}(p)=_0^{\mathrm{}}\frac{\rho (s)ds}{s+p^2}.$$
This expression can of course be obtained by taking the Fourier transform of both sides of Eq. (20).
With the explicit form of the propagator in configuration space given by Eq. (7), the representation in Eq. (20) turns out to be a particular example of the Hankel transform, namely the $`K`$-transform . Up to inessential factors of $`x`$ and $`m`$, Eq. (20) reduces to the generic form of the $`K`$-transform for a conjugate pair of functions $`f`$ and $`g`$,
$$g(x)=_0^{\mathrm{}}f(y)K_\nu (xy)\sqrt{xy}๐y.$$
(21)
The inverse of this transform is known to be given by
$$f(y)=\frac{1}{\pi i}_{ci\mathrm{}}^{c+i\mathrm{}}g(x)I_\nu (xy)\sqrt{xy}๐x$$
(22)
where $`I_\nu (x)`$ is a modified Bessel function of the first kind and the integration runs along a vertical contour in the complex plane to the right of the right-most singularity of the function $`g(x)`$ . In order to obtain a representation for the spectral density $`\rho (m^2)`$ of a sunrise-type diagram in general $`D`$-dimensional space-time one needs to apply the inverse $`K`$-transform to the particular case given by Eq. (20). One has
$$m^\lambda \rho (m^2)=\frac{(2\pi )^\lambda }{i}_{ci\mathrm{}}^{c+i\mathrm{}}\mathrm{\Pi }(x)x^{\lambda +1}I_\lambda (mx)๐x.$$
(23)
From Eq. (23) we obtain an explicit analytical representation for the spectral density $`\rho (s)`$ as a contour integral of the polarization function which reads
$$\rho (s)=\frac{(2\pi )^\lambda }{is^{\lambda /2}}_{ci\mathrm{}}^{c+i\mathrm{}}I_\lambda (x\sqrt{s})\mathrm{\Pi }(x)x^{\lambda +1}๐x$$
(24)
where $`I_\lambda (z)`$ is a modified Bessel function of the first kind and the integration runs along a vertical contour in the complex plane to the right of the right-most singularity of $`\mathrm{\Pi }(x)`$.
The inverse transform given by Eq. (24) completely solves the problem of determining the spectral density $`\rho (s)`$ of the general class of sunrise-type diagrams by reducing it to the computation of a one-dimensional integral along the contour in a complex plane. This is valid for the general class of sunrise-type diagrams with any number of internal lines and different masses. For $`(n+1)`$ internal lines with $`(n+1)`$ equal masses $`m`$ the spectral density reads
$$\rho (s)=\frac{m^{\lambda (n+1)}}{i(2\pi )^{n\lambda +n+1}s^{\lambda /2}}_{ci\mathrm{}}^{c+i\mathrm{}}I_\lambda (x\sqrt{s})\left(K_\lambda (mx)\right)^{n+1}x^{1n\lambda }๐x.$$
(25)
Because the contour can bypass the area of small values of $`x`$, the integral is finite, as it was stated before. The spectral density in turn can be used to restore the finite part of the correlation function in the momentum space representation. The path traced here is is an alternative to the calculation of the Fourier transform for sunrise-type diagrams.
In the standard, or momentum representation, the polarization function $`\stackrel{~}{\mathrm{\Pi }}(p)`$ is calculated from a $`n`$-loop diagram with $`n`$ $`D`$-dimensional integrations over the entangled loop momenta. It is clear that the momentum space evaluation becomes very difficult when the number of internal lines becomes large.
In order to demonstrate the applicability of our method described in more detail in , we show in Fig. 2 the result for the spectral density for the four-line sunrise-type diagram with $`D=e=2.718\mathrm{}`$, $`D=3`$, and $`D=\pi =3.14\mathrm{}`$. We have chosen these exotic values for $`D`$ to demonstrate the power of configuration space techniques. To the best of our knowledge there is no other method which allows one to compute the spectral density with such completeness and ease, i.e. through a one-dimensional integral. The technique was first suggested in and has been studied in detail in Refs. .
Note, finally, that the spectral density of a sunrise-type diagram is a representation of multi-particle phase space with the number of particles equal to the number of lines of the diagram. Our formula for the spectral density is a striking example of how the configuration-space technique can be used to determine the phase space of physical processes. Efficient techniques for the calculation of the multi-particle phase space are important for Monte-Carlo simulations and cross section calculations for multi-particle final states as they will occur at the upcoming Large Hadron Collider (LHC) under construction at CERN.
## 3 โTool boxโ for a practical analysis
In this section we give examples and practical prescriptions of how to proceed with the calculation of sunrise-type diagrams containing only a single Bessel function and powers of $`x`$. The massless case can be solved analytically in closed form for any number $`(n+1)`$ of internal lines since it contains only a single Bessel function and powers of $`x`$. The case for one massive line and $`n`$ massless lines (with two Bessel functions) is completely solvable analytically as well. The first nontrivial situation emerges with two massive lines. In this case a complete analytical solution is unknown in most of the cases, the exception being the case of equal masses that can be integrated analytical through hypergeometric functions. However, when one separates the diagram into a singular or pole part and a finite part, both parts can be calculated either analytically (for the pole part) or numerically (for the finite part) in very general situations. This will be demonstrated in this section.
### 3.1 Extraction of poles in dimensional regularization
Renormalization counterterms can be constructed analytically for general mass configurations and any number of loops. As mentioned earlier, in dimensional regularization the UV divergences become manifest as poles in $`\epsilon DD_0`$. In configuration space the UV divergences are related to short distances. In order to analyze the structure of singularities at small $`x`$, one has to expand the massive propagators at small $`x`$ (or, effectively, in the masses). In order to obtain the requisite conterterms one then integrates over $`x`$ which is analytically feasible for $`DD_0`$, i.e. $`\epsilon 0`$. Nevertheless, one should ensure convergence also at large $`x`$ which requires an appropriate IR regularization. Only this guarantees that singularities that emerge after integration are related to small $`x`$ divergences only (see e.g. ). This technical constraint can be met for instance by keeping one (or two) of the massive propagators in the integrand unexpanded. Since the massive propagator falls off exponentially at large $`x`$, this procedure is sufficient to ensure convergence at large $`x`$.
The possibility of retaining one propagator unexpanded exists because the integrals of one or a product of two Bessel functions are known analytically . This nice feature provides a tool for obtaining counterterms analytically. The necessary formulae read
$`{\displaystyle _0^{\mathrm{}}}x^{\mu 1}K_\nu (mx)๐x`$ $`=`$ $`2^{\mu 2}m^\mu \mathrm{\Gamma }\left({\displaystyle \frac{\mu +\nu }{2}}\right)\mathrm{\Gamma }\left({\displaystyle \frac{\mu \nu }{2}}\right)`$ (26)
$`{\displaystyle _0^{\mathrm{}}}x^{2\alpha 1}K_\mu (mx)K_\mu (mx)๐x`$ $`=`$ $`{\displaystyle \frac{2^{2\alpha 3}}{m^{2\alpha }\mathrm{\Gamma }(2\alpha )}}\mathrm{\Gamma }(\alpha +\mu )\mathrm{\Gamma }(\alpha )\mathrm{\Gamma }(\alpha )\mathrm{\Gamma }(\alpha \mu ).`$ (27)
In general, the finite parts of the relevant diagrams resulting from the subtraction of counterterms can be calculated only numerically. Still there are some simple examples where an analytic calculation is feasible. One of these examples is the three-loop sunrise-type diagram with two massive and two massless lines at vanishing external momentum, the three-loop vacuum bubble. There exist an analytical expression for the value of this particular vacuum bubble in configuration space. It is given by
$$\stackrel{~}{\mathrm{\Pi }}(0)=D(x,m)^2D(x,0)^2d^Dx=\left(\frac{(mx)^\lambda K_\lambda (mx)}{(2\pi )^{\lambda +1}x^{2\lambda }}\right)^2\left(\frac{\mathrm{\Gamma }(\lambda )}{4\pi ^{\lambda +1}x^{2\lambda }}\right)^2d^Dx.$$
(28)
While the angular integration in $`D`$-dimensional space-time is trivial, the problem of the residual radial integration is solved by using Eq. (27). The result for the integral in Eq. (28) reads
$$\stackrel{~}{\mathrm{\Pi }}(0)=\frac{(m^2)^{3\lambda 1}}{(4\pi )^{3(\lambda +1)}}\frac{\mathrm{\Gamma }(\lambda )^2\mathrm{\Gamma }(1\lambda )\mathrm{\Gamma }(12\lambda )^2\mathrm{\Gamma }(13\lambda )}{\mathrm{\Gamma }(\lambda +1)\mathrm{\Gamma }(24\lambda )}.$$
(29)
This result corresponds to the quantity $`M_1`$ in Ref. where it constitutes the simplest basis element for the computation of massive three-loop diagrams in a general three-loop topology. Again, any number of massless lines can be added.
The result given by Eq. (29) is obtained in a concise form and valid for any dimension $`D`$. Therefore its pole parts can be explicitly singled out through a direct expansion into a Laurent series near an integer value of the space-time dimension $`D`$. Taking $`D=42\epsilon `$ (i.e. $`\lambda =1\epsilon `$) and expanding the result for small $`\epsilon `$ one readily finds
$$\stackrel{~}{\mathrm{\Pi }}(0)=\frac{(m^2)^{23\epsilon }}{3(4\pi )^{3(2\epsilon )}}\left(\frac{1}{\epsilon ^3}+\frac{7}{2\epsilon ^2}+\frac{25+6\zeta (2)}{4\epsilon }\frac{542\zeta (2)56\zeta (3)}{8}+O(\epsilon )\right).$$
Using this example one can take a closer look at the two-dimensional case. As was stated before, the case $`D_0=2`$ is much softer and the UV divergences are practically absent. Hence one could come to the conclusion that the expression given by Eq. (29) should be perfectly finite at $`D=D_0=2`$. In terms of $`D=2\lambda +2`$ the limit corresponds to taking $`\lambda 0`$. Therefore, one might be surprised to find a singularity $`\mathrm{\Gamma }(\lambda )^2`$ in Eq. (29). However, this is not a UV singularity. Because massless propagators are not defined in the two-dimensional world, the limit $`\lambda 0`$ is ill-defined in this case. Indeed, looking at Eq. (8), the singularity of this sunrise-type diagram is exactly the product of two singular massless propagators. Power counting for the integral expression in Eq. (8) shows that the singularity is caused by the infrared (IR) region. As this is an IR problem, the afore mentioned BPHZ theorem does not apply and the divergence cannot be treated with local counterterms in the usual manner.
Returning to the more general case, a closed-form evaluation (as in Eq. (29)) is not possible and an explicit subtraction is required to separate singular and finite parts. In our case it is convenient to use momentum subtraction which in fact is the oldest renormalization method. Momentum subtraction consists of subtracting a polynomial at some fixed momentum point. For massive diagrams the expansion point $`p=0`$ is safe (i.e. has no infrared (IR) singularity), and the prescription is realized by expanding the function
$$\left(\frac{px}{2}\right)^\lambda J_\lambda (px)$$
(30)
(which is the kernel or weight function of the integral transformation in Eq. (18)) in a Taylor series around $`p=0`$ in terms of a polynomial series in $`p^2`$ (cf. Eq. (A4) in Appendix A). The subtraction at order $`N`$ is achieved by writing
$$\left[\left(\frac{px}{2}\right)^\lambda J_\lambda (px)\right]_N=\left(\frac{px}{2}\right)^\lambda J_\lambda (px)\underset{k=0}{\overset{N}{}}\frac{(1)^k}{k!\mathrm{\Gamma }(\lambda +k+1)}\left(\frac{px}{2}\right)^{2k}$$
(31)
and by keeping $`N`$ terms in the expansion on the right hand side. After substituting the expansion Eq. (31) into Eq. (18) one obtains the momentum subtracted polarization function
$$\stackrel{~}{\mathrm{\Pi }}_{\mathrm{mom}}(p)=\stackrel{~}{\mathrm{\Pi }}(p)\underset{k=0}{\overset{N}{}}\frac{p^{2k}}{k!}\left(\frac{d}{dp^2}\right)^k\stackrel{~}{\mathrm{\Pi }}(p)|_{p^2=0}$$
(32)
which is finite if the number of subtractions $`N`$ is sufficiently high. The function $`\stackrel{~}{\mathrm{\Pi }}(p)`$ is divergent as well as any of the derivatives on the right hand side of Eq. (32). The divergences require regularization. However, the difference, i.e. the quantity $`\stackrel{~}{\mathrm{\Pi }}_{\mathrm{mom}}(p)`$ is finite and independent of any regularization used to give a meaning to each individual term in Eq. (32). Note that the expansion in Eq. (31) is a polynomial in $`p^2`$ in accordance with the general structure of the $`R`$-operation . The number $`N`$ of necessary subtractions is determined by the divergence index of the diagram and can be found according to the standard power counting rules. The subtraction at the origin $`p=0`$ is allowed if there is at least one massive line in the diagram along with an arbitrary number of massless lines. If there are no massive internal lines at all, the corresponding diagram can easily be calculated analytically and the problem of subtraction is trivial. After having performed the requisite subtraction, one can take the limit $`DD_0`$ in Eq. (18) where $`D_0`$ is an integer. The diagram as a whole becomes finite after subtraction which reflects the topology of the sunrise-type diagram: there is no divergent subdiagram in the sense of Bogoliubov.
In selecting the momentum subtraction in order to regularize the singularities we have selected a particular subtraction or renormalization scheme. The coefficients $`C_i`$ in Eq. (19) depend on this choice. While the highest coefficient related to the strongest singularity is unique, the other coefficients are scheme dependent. The standard scheme used in the literature is the $`\overline{\mathrm{MS}}`$ scheme. The transition from the momentum subtraction scheme to the $`\overline{\mathrm{MS}}`$ scheme is achieved by the familiar redefinition of the measure
$$d^Dx(4\pi \mu ^2e^{\gamma _E})^\epsilon d^Dx$$
(33)
with $`\mu `$ being the renormalization scale within dimensional regularization and $`\gamma _E`$ being Eulerโs constant $`\gamma _E=0.577\mathrm{}`$
### 3.2 Singular and finite parts for particular parameter values of the genuine sunrise: warm-up example demonstrating the technique
As an example of calculating the counterterms and evaluating the remaining finite part of the diagram we present the special case of the genuine sunrise diagram with three massive lines, i.e. the two-loop sunrise-type diagram. For special values of the external momenta a closed-form expression for the genuine sunrise diagram is known within dimensional regularization using momentum space techniques . We shall reproduce this result in the following. The one-dimensional integral to be analyzed is given by
$$\stackrel{~}{\mathrm{\Pi }}(p)=2\pi ^{\lambda +1}_0^{\mathrm{}}\left(\frac{px}{2}\right)^\lambda J_\lambda (px)D(x,m_1)D(x,m_2)D(x,m_3)x^{2\lambda +1}๐x.$$
(34)
To determine the finite part we first use momentum subtraction and split $`\stackrel{~}{\mathrm{\Pi }}(p)`$ up into its finite and infinite (but dimensionally regularized) part,
$$\stackrel{~}{\mathrm{\Pi }}(p)=\stackrel{~}{\mathrm{\Pi }}_{\mathrm{mom}}(p)+\stackrel{~}{\mathrm{\Pi }}_{\mathrm{sing}}(p)$$
(35)
where $`\stackrel{~}{\mathrm{\Pi }}_{\mathrm{mom}}(p)`$ is a momentum subtracted polarization function and $`\stackrel{~}{\mathrm{\Pi }}_{\mathrm{sing}}(p)`$ is a counterterm in dimensional regularization. Power counting shows that only two subtractions are necessary. The explicit expression for the momentum subtracted polarization function reads
$`\stackrel{~}{\mathrm{\Pi }}_{\mathrm{mom}}(p)`$ $`=`$ $`2\pi ^{\lambda +1}{\displaystyle _0^{\mathrm{}}}\left[\left({\displaystyle \frac{px}{2}}\right)^\lambda J_\lambda (px)\right]_1D(x,m_1)D(x,m_2)D(x,m_3)x^{2\lambda +1}๐x`$ (36)
$`=`$ $`2\pi ^{\lambda +1}{\displaystyle _0^{\mathrm{}}}\left[\left({\displaystyle \frac{px}{2}}\right)^\lambda J_\lambda (px){\displaystyle \frac{1}{\mathrm{\Gamma }(\lambda +1)}}+{\displaystyle \frac{p^2x^2}{4}}{\displaystyle \frac{1}{\mathrm{\Gamma }(\lambda +2)}}\right]`$
$`\times D(x,m_1)D(x,m_2)D(x,m_3)x^{2\lambda +1}dx`$
while the pole part is given by a first order polynomial in $`p^2`$ of the form $`A+Bp^2`$. The actual values of the coefficients $`A`$ and $`B`$ depend on the regularization scheme that has been used. We use dimensional regularization and obtain
$`\stackrel{~}{\mathrm{\Pi }}_{\mathrm{sing}}(p)`$ $`=`$ $`A+p^2B`$ (37)
$`=`$ $`{\displaystyle \frac{2\pi ^{\lambda +1}}{\mathrm{\Gamma }(\lambda +1)}}{\displaystyle _0^{\mathrm{}}}D(x,m_1)D(x,m_2)D(x,m_3)x^{2\lambda +1}๐x`$
$`p^2{\displaystyle \frac{2\pi ^{\lambda +1}}{4\mathrm{\Gamma }(\lambda +2)}}{\displaystyle _0^{\mathrm{}}}x^2D(x,m_1)D(x,m_2)D(x,m_3)x^{2\lambda +1}๐x.`$
With the help of the representation given in Eq. (37) we can apply our strategy in a straightforward manner. The coefficients $`A`$ and $`B`$ are simple numbers independent of $`p`$. They contain divergent parts (regularized within dimensional regularization) and need to be computed only once to recover the function $`\stackrel{~}{\mathrm{\Pi }}(p)`$ for any $`p`$. In the momentum subtracted part one can forego the regularization (it is finite by the $`R`$-operation) and perform the one-dimensional integration numerically for $`D=4`$ (i.e. $`\lambda =1`$).
In the particular case of the genuine sunrise diagram the necessary integrals are known analytically in closed form and can be found in integral tables (see e.g. ). However, using these tables may not always be convenient and we therefore again present a simplified approach which allows one to deal with even more complicated cases in a simpler manner. Let us specify to the particular case $`p=m_1+m_2m_3`$ (which corresponds to a pseudothreshold) where an analytical answer for the integral exists . For simplicity we choose $`m_1=m_2=m_3/2=m`$. Then $`p=0`$ and $`\stackrel{~}{\mathrm{\Pi }}_{\mathrm{mom}}(0)=0`$ (which is a regular function at this Euclidean point). In this case the counterterm $`Bp^2`$ vanishes because the quantity $`B`$ is finite at finite $`\epsilon `$. We thus only have to consider the coefficient $`A`$. On the other hand our considerations are completely general since for arbitrary $`p`$ one only requires more terms in the $`p^2`$-expansion. For the special mass configuration considered here the result of reads
$$\stackrel{~}{\mathrm{\Pi }}^{\mathrm{ref}}(0)=\pi ^{42\epsilon }\frac{m^{24\epsilon }\mathrm{\Gamma }^2(1+\epsilon )}{(1\epsilon )(12\epsilon )}\left[\frac{3}{\epsilon ^2}+\frac{8\mathrm{ln}2}{\epsilon }8\mathrm{ln}^22\right]+O(\epsilon ).$$
(38)
This result can be extracted from the first line of Eq. (37) using the integral tables given in . However, in the general mass case the necessary formulas are rather cumbersome. Even for the special mass configuration considered here they are not so simple. We therefore discuss a short cut which allows one to obtain results immediately without having to resort to integral tables. What we really need is an expansion in $`\epsilon `$. Basically we need the integral
$$_0^{\mathrm{}}D(x,m)D(x,m)D(x,2m)x^{2\lambda +1}๐x$$
(39)
which is of the general form
$$_0^{\mathrm{}}x^\rho K_\mu (mx)K_\mu (mx)K_\mu (2mx)๐x$$
(40)
($`\mu =\lambda `$ and $`\rho =1\lambda `$ in our case). For a product of two Bessel functions in the integrand (without the last one in the above equation) the result of the integration is known analytically and is given by Eq. (27). Let us reduce the problem at hand to Eq. (27) and do our numerical evaluations with functions in four-dimensional space-time where no regularization is necessary. To do so we subtract the leading singularities at small $`\xi `$ from the last Bessel function in Eq. (40) using the series expansion near the origin given by Eq. (A10). After decomposing the whole answer into finite and singular parts according to $`(2\pi )^{2D}A=F+S`$ (where the total normalization of has been adopted in the definition of $`F`$ and $`S`$) we find for the singular part
$`S`$ $`=`$ $`{\displaystyle \frac{(2\pi )^Dm^{2\lambda }}{\mathrm{\Gamma }(\lambda +1)}}{\displaystyle _0^{\mathrm{}}}x^{2(1\lambda )1}K_\lambda (mx)K_\lambda (mx){\displaystyle \frac{\mathrm{\Gamma }(\lambda )}{2}}\left[1+{\displaystyle \frac{(mx)^2}{1\lambda }}{\displaystyle \frac{\mathrm{\Gamma }(1\lambda )}{\mathrm{\Gamma }(1+\lambda )}}(mx)^{2\lambda }\right]๐x`$ (41)
$`=`$ $`{\displaystyle \frac{(2\pi )^Dm^{24\epsilon }}{\mathrm{\Gamma }(\lambda +1)}}{\displaystyle _0^{\mathrm{}}}\xi ^{2\epsilon 1}K_\lambda (\xi )K_\lambda (\xi ){\displaystyle \frac{\mathrm{\Gamma }(\lambda )}{2}}\left[1+{\displaystyle \frac{\xi ^2}{1\lambda }}{\displaystyle \frac{\mathrm{\Gamma }(1\lambda )}{\mathrm{\Gamma }(1+\lambda )}}\xi ^{2\lambda }\right]๐\xi `$
$`=`$ $`\pi ^{42\epsilon }{\displaystyle \frac{m^{24\epsilon }\mathrm{\Gamma }^2(1+\epsilon )}{(1\epsilon )(12\epsilon )}}\left[{\displaystyle \frac{3}{\epsilon ^2}}+{\displaystyle \frac{8\mathrm{ln}2}{\epsilon }}+8(22\mathrm{ln}2\mathrm{ln}^22)\right]+O(\epsilon ).`$
The pole contributions coincide with the result in Eq. (38) while the finite part is different. It is corrected by the finite expression
$`F`$ $`=`$ $`{\displaystyle \frac{(2\pi )^Dm^{2\lambda }}{\mathrm{\Gamma }(\lambda +1)}}{\displaystyle _0^{\mathrm{}}}x^{2(1\lambda )1}K_\lambda (mx)K_\lambda (mx)`$ (42)
$`\times \left\{(mx)^\lambda K_\lambda (2mx){\displaystyle \frac{\mathrm{\Gamma }(\lambda )}{2}}\left[1+{\displaystyle \frac{(mx)^2}{1\lambda }}{\displaystyle \frac{\mathrm{\Gamma }(1\lambda )}{\mathrm{\Gamma }(1+\lambda )}}(mx)^{2\lambda }\right]\right\}dx.`$
Because $`F`$ is finite (no strong singularity at small $`x`$) one can put $`\lambda =1`$ to obtain
$$F=16\pi ^4m^2_0^{\mathrm{}}\frac{dx}{x}K_1(x)K_1(x)\left\{xK_1(2x)\frac{1}{2}\left[1+x^2(1+2\gamma _E+2\mathrm{ln}x)\right]\right\}.$$
(43)
Doing the numerical integration in Eq. (43) results in the expression
$$F=16\pi ^4m^2[0.306853\mathrm{}].$$
(44)
One can restore all $`\epsilon `$-dependence in the normalization factors as in Eqs. (38) and (41) because $`F`$ is not singular in $`\epsilon `$, so that this change of normalization is absorbed by the $`O(\epsilon )`$ symbol. One obtains
$$F=\pi ^{42\epsilon }\frac{m^{24\epsilon }\mathrm{\Gamma }^2(1+\epsilon )}{(1\epsilon )(12\epsilon )}\left[0.306853\mathrm{}\right]+O(\epsilon ).$$
(45)
Adding $`F`$ from Eq. (45) to $`S`$ from Eq. (41) one obtains a numerical result for the finite part of Eq. (38). In order to establish the full coincidence with the result of Eq. (38), note that
$$0.306853\mathrm{}=(1\mathrm{ln}2)\times 1.0000\mathrm{}.$$
Of course we have used the analytical expression $`(1\mathrm{ln}2)`$ for the numerical quantity $`0.306853\mathrm{}`$ for illustrative reasons because we knew the final answer.
We emphasize that there is nothing new in computing the polarization function related to this diagram at any $`p`$. Some finite part appears from the momentum subtracted polarization function. One then needs another counterterm for non-zero $`p^2`$. The computation of these additional terms proceeds in analogy to what has been done before. It is even simpler because the singularity at small $`x`$ is weaker and only one subtraction of the Bessel function is necessary. This computational technique works for any complex value $`p^2`$.
### 3.3 Three-loop sunrise-type diagram for special momenta
In a similar manner one can treat the three-loop sunrise-type diagram with two equal masses $`M`$, one mass $`m`$ with values between zero and $`M`$ and a massless line for outer momentum $`p^2=m^2`$. This particular sunrise-type diagram was analyzed by Mastrolia whose results we confirm. The starting integral expression reads
$`\stackrel{~}{\mathrm{\Pi }}(p)`$ $`=`$ $`2\pi ^{\lambda +1}{\displaystyle _0^{\mathrm{}}}\left({\displaystyle \frac{px}{2}}\right)^\lambda J_\lambda (px)D(x,M)D(x,M)D(x,m)D(x,0)x^{2\lambda +1}๐x`$ (46)
$`=`$ $`{\displaystyle \frac{M^{2\lambda }m^\lambda \mathrm{\Gamma }(\lambda )}{2(2\pi )^{3(\lambda +1)}}}{\displaystyle _0^{\mathrm{}}}\left({\displaystyle \frac{px}{2}}\right)^\lambda J_\lambda (px)K_\lambda (Mx)K_\lambda (Mx)K_\lambda (mx)x^{13\lambda }๐x.`$
The integral has the same form as in the previous example. It is obvious that we can proceed in the same way as before: first expand the integrand in $`p`$ and $`m`$ in order to find the counterterms, and second determine the finite contribution for $`p^2=m^2`$ numerically. However, in this case we have to expand the Bessel function $`J_\lambda (px)`$ up to the second order in $`p^2`$. For the singular part we obtain
$`\stackrel{~}{\mathrm{\Pi }}_{\mathrm{sing}}(p)`$ $`=`$ $`A+p^2B+{\displaystyle \frac{1}{2}}p^4C`$ (47)
$`=`$ $`{\displaystyle \frac{2\pi ^{\lambda +1}}{\mathrm{\Gamma }(\lambda +1)}}{\displaystyle _0^{\mathrm{}}}D(x,m_1)D(x,m_2)D(x,m_3)x^{2\lambda +1}๐x`$
$`p^2{\displaystyle \frac{2\pi ^{\lambda +1}}{4\mathrm{\Gamma }(\lambda +2)}}{\displaystyle _0^{\mathrm{}}}x^2D(x,m_1)D(x,m_2)D(x,m_3)x^{2\lambda +1}๐x`$
$`+p^4{\displaystyle \frac{2\pi ^{\lambda +1}}{32\mathrm{\Gamma }(\lambda +3)}}{\displaystyle _0^{\mathrm{}}}x^4D(x,m_1)D(x,m_2)D(x,m_3)x^{2\lambda +1}๐x`$
where
$`A`$ $`=`$ $`{\displaystyle \frac{2M^{46\epsilon }}{(4\pi )^{3(2\epsilon )}\mathrm{\Gamma }(22\epsilon )\mathrm{\Gamma }(2\epsilon )}}({\displaystyle \frac{1+2y^2}{6\epsilon ^3}}+{\displaystyle \frac{2+8y^2y^424y^2\mathrm{ln}y}{24\epsilon ^2}}`$
$`+{\displaystyle \frac{1}{\epsilon }}({\displaystyle \frac{3}{16}}(2+y^4)+{\displaystyle \frac{2}{3}}(1+2y^2)\zeta (2)y^2\mathrm{ln}y+{\displaystyle \frac{1}{4}}y^4\mathrm{ln}y+y^2\mathrm{ln}^2y)),`$
$`B`$ $`=`$ $`{\displaystyle \frac{2M^{26\epsilon }}{(4\pi )^{3(2\epsilon )}\mathrm{\Gamma }(22\epsilon )\mathrm{\Gamma }(3\epsilon )}}\left({\displaystyle \frac{2+y^2}{12\epsilon ^2}}+{\displaystyle \frac{2+y^26y^2\mathrm{ln}y}{12\epsilon }}\right),`$
$`C`$ $`=`$ $`{\displaystyle \frac{2M^{6\epsilon }}{(4\pi )^{3(2\epsilon )}\mathrm{\Gamma }(22\epsilon )\mathrm{\Gamma }(4\epsilon )}}\left({\displaystyle \frac{1}{6\epsilon }}\right)`$ (48)
and $`y=m/M`$.
### 3.4 Infinite parts of sunrise-type diagrams with different masses at any number of loops: general techniques
If one deals with sunrise-type diagrams with different masses one has to be careful to use a symmetric subtraction procedure. Otherwise one looses the symmetry of the original diagram. Such an asymmetry would arise if one would expand one of the Bessel functions in the integrand of the sunrise-type diagram as was done in the previous example for the massive line with mass $`m`$. If one wants to obtain a result which is symmetric under the exchange of the masses $`m_i`$ as expected from the topology of the diagram and from the initial form of the integral we have to subtract a counterterm which is symmetric in these masses. In calculating the genuine sunrise diagram one can follow this procedure by using a damping function such as $`e^{\mu ^2x^2}`$ to suppress contributions for large values of $`x`$, corresponding to IR singularities. This is done by introducing the factor $`1=e^{\mu ^2x^2}e^{\mu ^2x^2}`$ into the integrand. In this case all massive propagators along with the factor $`e^{\mu ^2x^2}`$ can be expanded in terms of small $`x`$ values. The expansion of the propagators is straightforward and can be obtained from Eqs. (A4) and (A10) in Appendix A. The final integration can then be done formally by using the identity
$$_0^{\mathrm{}}x^{p1}e^{\mu ^2x^2}๐x=\frac{1}{2}\mu ^p\mathrm{\Gamma }(p/2)$$
(49)
which is nothing but the definition of Eulerโs Gamma function.
As an illustration we compute the four-loop bubble diagram which was also calculated by Laporta whose results we confirm. In the sample results below we extract the normalization factor
$$๐ฉ_n=\left(\frac{(4\pi )^{2\epsilon }}{\mathrm{\Gamma }(1+\epsilon )}\right)^n$$
(50)
for the $`n`$-loop (or $`(n+1)`$-line) sunrise-type diagram. Extraction of this factor renders all pole parts to be rational numbers. The results of the calculation are
$`๐ฉ_1\stackrel{~}{\mathrm{\Pi }}_1(p^2)`$ $`=`$ $`{\displaystyle \frac{1}{\epsilon }}+O(\epsilon ^0),`$
$`๐ฉ_2\stackrel{~}{\mathrm{\Pi }}_2(p^2)`$ $`=`$ $`m^2\left\{{\displaystyle \frac{3}{2\epsilon ^2}}{\displaystyle \frac{9}{2\epsilon }}\right\}{\displaystyle \frac{p^2}{4\epsilon }}+O(\epsilon ^0),`$
$`๐ฉ_3\stackrel{~}{\mathrm{\Pi }}_3(p^2=0)`$ $`=`$ $`m^4\left\{{\displaystyle \frac{2}{\epsilon ^3}}+{\displaystyle \frac{23}{3\epsilon ^2}}+{\displaystyle \frac{35}{2\epsilon }}\right\}+O(\epsilon ^0),`$
$`๐ฉ_3\stackrel{~}{\mathrm{\Pi }}_3(p^2=m^2)`$ $`=`$ $`m^4\left\{{\displaystyle \frac{2}{\epsilon ^3}}+{\displaystyle \frac{22}{3\epsilon ^2}}+{\displaystyle \frac{577}{36\epsilon }}\right\}+O(\epsilon ^0),`$
$`๐ฉ_4\stackrel{~}{\mathrm{\Pi }}_4(p^2=0)`$ $`=`$ $`m^6\left\{{\displaystyle \frac{5}{2\epsilon ^4}}{\displaystyle \frac{35}{3\epsilon ^3}}{\displaystyle \frac{4565}{144\epsilon ^2}}{\displaystyle \frac{58345}{864\epsilon }}\right\}+O(\epsilon ^0),`$
$`๐ฉ_4\stackrel{~}{\mathrm{\Pi }}_4(p^2=m^2)`$ $`=`$ $`m^6\left\{{\displaystyle \frac{5}{2\epsilon ^4}}{\displaystyle \frac{45}{4\epsilon ^3}}{\displaystyle \frac{4255}{144\epsilon ^2}}{\displaystyle \frac{106147}{1728\epsilon }}\right\}+O(\epsilon ^0),`$
$`๐ฉ_5\stackrel{~}{\mathrm{\Pi }}_5(p^2=0)`$ $`=`$ $`m^8\left\{{\displaystyle \frac{3}{\epsilon ^5}}+{\displaystyle \frac{33}{2\epsilon ^4}}+{\displaystyle \frac{1247}{24\epsilon ^3}}+{\displaystyle \frac{180967}{1440\epsilon ^2}}+{\displaystyle \frac{898517}{3456\epsilon }}\right\}+O(\epsilon ^0),`$
$`๐ฉ_5\stackrel{~}{\mathrm{\Pi }}_5(p^2=m^2)`$ $`=`$ $`m^8\left\{{\displaystyle \frac{3}{\epsilon ^5}}+{\displaystyle \frac{16}{\epsilon ^4}}+{\displaystyle \frac{49}{\epsilon ^3}}+{\displaystyle \frac{6967}{60\epsilon ^2}}+{\displaystyle \frac{1706063}{7200\epsilon }}\right\}+O(\epsilon ^0),`$
$`๐ฉ_6\stackrel{~}{\mathrm{\Pi }}_6(p^2=0)`$ $`=`$ $`m^{10}\{{\displaystyle \frac{7}{2\epsilon ^6}}{\displaystyle \frac{133}{6\epsilon ^5}}{\displaystyle \frac{238}{3\epsilon ^4}}{\displaystyle \frac{77329}{360\epsilon ^3}}`$ (51)
$`{\displaystyle \frac{21221921}{43200\epsilon ^2}}{\displaystyle \frac{2596372387}{2592000\epsilon }}\}+O(\epsilon ^0).`$
The coefficient of the leading singularity in $`\epsilon `$ is independent of $`p^2`$. In Appendix D we list results for unequal mass configurations up to four-loop order. When setting all masses equal the results of the general mass case can be seen to agree with the above equal mass results.
### 3.5 Finite part for the genuine sunrise with three different masses
In configuration space the genuine sunrise diagram with three different masses $`m_1`$, $`m_2`$, and $`m_3`$ is given by
$$\mathrm{\Pi }(x)=D(x,m_1)D(x,m_2)D(x,m_3).$$
(52)
The Fourier transform of this polarization function in the Euclidean domain reads
$$\stackrel{~}{\mathrm{\Pi }}(p)=2\pi ^{\lambda +1}_0^{\mathrm{}}\left(\frac{px}{2}\right)^\lambda J_\lambda (px)D(x,m_1)D(x,m_2)D(x,m_3)x^{2\lambda +1}๐x.$$
(53)
In dimensional regularization the singular parts ($`D=42\epsilon `$, i.e. $`\lambda =1\epsilon `$) are given by
$$\stackrel{~}{\mathrm{\Pi }}_{\mathrm{sing}}(p)=\frac{\mu ^{4\epsilon }}{\pi ^{42\epsilon }}\left\{\frac{m_1^2+m_2^2+m_3^2}{512\epsilon ^2}+\frac{1}{\epsilon }\left(\underset{i=1}{\overset{3}{}}\frac{m_i^2\mathrm{ln}(m_ie^\gamma /2\mu )}{128}\frac{p^2}{1024}\right)\right\}$$
(54)
where $`\gamma ^{}=\gamma _E/23/4`$. The counterterms for the Fourier transform are obtained as explained before by introducing $`1=e^{\mu ^2x^2}e^{\mu ^2x^2}`$ and expanding the relevant Bessel functions together with $`e^{\mu ^2x^2}`$. This is the reason for the explicit dependence on the arbitrary parameter $`\mu `$. In this way one readily finds the pole parts for any diagram with any mass configuration, i.e. one finds the coefficients $`C_i`$ of the counterterms in Eq. (19). Even if an external momentum goes through the diagram, the pole parts can be still calculated easily since they are given by integrals over the expanded integrand. The momentum subtracted finite part is given by
$$\stackrel{~}{\mathrm{\Pi }}_{\mathrm{fin}}(p)=2\pi ^{\lambda +1}_0^{\mathrm{}}\left[\left(\frac{px}{2}\right)^\lambda J_\lambda (px)\right]_ND(x,m_1)D(x,m_2)D(x,m_3)x^{2\lambda +1}e^{\mu ^2x^2}e^{\mu ^2x^2}๐x$$
(55)
where $`N`$ is the order of the expansion needed to ensure integrability. Naturally $`N`$ has the same order as the one used earlier in order to extract the counterterms. It is obvious, though, that the expansion can terminate at positive powers of $`x`$. Finally, one can set $`\epsilon =0`$ to calculate the finite part numerically. As a technical remark note that most numerical integration routines will run into problems if the upper limit extends to infinity. However, already for values of the order $`x1`$ the integrand is negligibly small, so that the integration can be cut off at this point. Also the region around the origin might cause trouble for the convergence of the numerical integration. In this case the integration interval can be subdivided with an interval close to the origin, and the whole integrand in this subinterval can be expanded in $`x`$ and integrated analytically. For specific values of $`p^2`$, $`m_1`$, $`m_2`$, and $`m_3`$ we were able to reproduce results given in the literature (see e.g. results given in Refs. ).
### 3.6 Calculation of the finite part: further examples
The examples presented so far are well-known and have been obtained before by using techniques different from the ones presented here. While we can numerically compute any sunrise-type diagram with any arbitrary number of internal massive lines it is not always possible to find corresponding analytical expressions to compare with. Beyond two loops only a few examples can be found in the literature. For instance, we have calculated the result for the vacuum bubble diagram shown in Fig. 3 (see also Sec. 5.4). For $`p0`$ we obtain (cf. Eq. (A4))
$$\left(\frac{px}{2}\right)^\lambda J_\lambda (px)\frac{1}{\mathrm{\Gamma }(\lambda +1)}.$$
(56)
One therefore has
$$\stackrel{~}{\mathrm{\Pi }}(0)=\frac{2\pi ^{\lambda +1}}{\mathrm{\Gamma }(\lambda +1)}_0^{\mathrm{}}\left(D^{(1)}(x,m)\right)^3D(x,m)x^{2\lambda +1}๐x$$
(57)
where $`D^{(1)}(x,m)`$ is defined in Eq. (12). In order to compare with the literature , we multiply with a relative factor $`๐ฉ_3m^{2+6\epsilon }`$ (cf. Eq. (50)) to obtain
$$B_\epsilon =๐ฉ_3m^{2+6\epsilon }\stackrel{~}{\mathrm{\Pi }}(0)=\frac{2^{22\epsilon }}{(1\epsilon )\mathrm{\Gamma }(1+\epsilon )^3\mathrm{\Gamma }(1\epsilon )}_0^{\mathrm{}}K_\epsilon ^3(x)K_{1\epsilon }(x)x^{2+2\epsilon }๐x.$$
(58)
The integral remains finite even for $`\epsilon 0`$ and, after finite integration, we obtain the numerical result
$$B_0=4_0^{\mathrm{}}K_0^3(x)K_1(x)x^2๐x=2.1035995805\mathrm{}$$
(59)
A comparison with the results presented in Ref. shows full numerical agreement. Since we know that the quantity $`B_0`$ contains only a single transcendental number, namely $`\zeta (3)`$, where $`\zeta (z)`$ is Riemannโs $`\zeta `$-function , we can take its numerical value $`\zeta (3)=1.202056903\mathrm{}`$ to obtain the result
$$B_0=\zeta (3)1.7500000\mathrm{}=\frac{7}{4}\zeta (3)$$
which is nothing but the analytical result given in Ref. .
We conclude that the configuration space method is simple and allows one to numerically compute diagrams with sunrise-type topology in a very efficient manner. When the structure of the transcendentality of the result is known for a particular diagram (as in the latest example where only $`\zeta (3)`$ is present) or when it can be obtained from an educated guess looking at the topology of the diagram our numerical technique can be used to restore the rational coefficients of these transcendentalities which can be thought of as elements of the basis for a class of diagrams.
### 3.7 Higher order expansion in $`\epsilon `$: <br>general techniques and the four-loop vacuum bubble
After having determined the singular parts of the Laurent series expansion using the damping factor method what remains to be done is to calculate the coefficients of the positive powers of $`\epsilon `$ in the $`\epsilon `$-expansion including the finite $`\epsilon ^0`$ term. Higher order terms in the $`\epsilon `$-expansion are needed if the sunrise-type diagram is inserted into a divergent diagram or when one is using the integration-by-parts recurrence relation which can generate inverse powers of $`\epsilon `$. In order to determine these higher order terms one needs to resort to the Bessel function method. What is technically needed is to develop a procedure for the $`\epsilon `$-expansion of Bessel functions.
The $`O(\epsilon )`$ term can be obtained analytically. Indeed, within dimensional regularization the propagator in configuration space contains Bessel functions with non-integer index depending on the regularization parameter $`\epsilon `$, as they occurred for instance in Sec. 3.6 at an intermediate step. In order to expand the Bessel function in the parameter $`\epsilon `$ entering its index, we use Eq. (A13) for the derivative of the modified Bessel function of the second kind with respect to its index near integer values of this index. For instance, we obtain
$`K_\epsilon (x)`$ $`=`$ $`K_0(x)+O(\epsilon ^2),`$
$`K_{1\epsilon }(x)`$ $`=`$ $`K_1(x){\displaystyle \frac{\epsilon }{x}}K_0(x)+O(\epsilon ^2),`$
$`K_{2\epsilon }(x)`$ $`=`$ $`K_2(x){\displaystyle \frac{2\epsilon }{x}}K_1(x){\displaystyle \frac{2\epsilon }{x^2}}K_0(x)+O(\epsilon ^2)\mathrm{}`$ (60)
For the more general cases one can use the integral representation
$$K_\nu (z)=_0^{\mathrm{}}\mathrm{cosh}(\nu t)e^{z\mathrm{cosh}t}๐t$$
(61)
for the index expansion. However, one can only obtain numerical results in this case. In expanding the integrand, we obtain
$`K_\epsilon (z)`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\epsilon ^{2n}}{(2n)!}}f_{2n}(z),`$
$`K_{1\epsilon }(z)`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\epsilon ^{2n}}{(2n)!}}a_{2n}(z){\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\epsilon ^{2n+1}}{(2n+1)!}}b_{2n+1}(z)`$ (62)
where
$`f_k(z)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}t^ke^{z\mathrm{cosh}t}๐t,`$
$`a_k(z)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}t^k\mathrm{cosh}te^{z\mathrm{cosh}t}dt,`$
$`b_k(z)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}t^k\mathrm{sinh}te^{z\mathrm{cosh}t}dt.`$ (63)
Due to the recurrence relations in the index $`\nu `$ for the family $`K_\nu (z)`$ (cf. Sec. 5.6) the two formulas in Eq. (3.7) suffice to calculate the Bessel function $`K_\nu (z)`$ for any $`\nu `$. The two functions $`a_k(z)`$ and $`b_k(z)`$ in Eq. (3.7) are in turn related to the functions $`f_k(z)`$,
$$a_k(z)=\frac{d}{dz}f_k(z),b_k(z)=\frac{k}{z}f_{k1}(z).$$
(64)
Using these relations, we obtain
$$K_{1\epsilon }(z)=\underset{n=0}{\overset{\mathrm{}}{}}\frac{\epsilon ^{2n}}{(2n)!}\left(\frac{d}{dz}+\frac{\epsilon }{z}\right)f_{2n}(z).$$
(65)
The family of functions $`f_k(z)`$ is rather close to the original set of Bessel functions $`K_\nu (z)`$ (cf. Fig. 4) and can easily be studied both analytically and numerically. The limits at $`z0`$ and $`z\mathrm{}`$ are known analytically and are simple. They allow for an efficient interpolation for intermediate values of the argument.
The functions $`f_k(z)`$ satisfy the differential equation
$$\left(\frac{d^2}{dz^2}+\frac{1}{z}\frac{d}{dz}1\right)f_k(z)=\frac{k(k1)}{z^2}f_{k2}(z)$$
(66)
related to the Bessel differential equation. The small $`z`$ behaviour
$$f_k(z)=\frac{1}{k+1}\mathrm{ln}^{k+1}\left(\frac{1}{z}\right)\left(1+O\left(\frac{1}{\mathrm{ln}(z)}\right)\right)$$
(67)
can be found by using yet another representation for the function $`f_k(z)`$,
$$f_k(z)=_1^{\mathrm{}}\frac{e^{zu}}{\sqrt{u^21}}\mathrm{ln}^k(u+\sqrt{u^21})๐u$$
(68)
or directly from the behaviour of the function $`K_\nu (x)`$ at small $`x`$,
$$K_\epsilon (x)=\frac{1}{\epsilon }\mathrm{sinh}\left(\epsilon \mathrm{ln}(1/x)\right).$$
(69)
While in four-dimensional space-time the massive propagator contains the Bessel function $`K_{1\epsilon }(z)`$, the basic function in $`D=2`$ dimensional space-time is the Bessel function $`K_\epsilon (z)`$ since the propagator reads
$$D(x,m)|_{D=2}=\frac{1}{2\pi }K_0(mx).$$
(70)
As an example for the numerical evaluation of the $`\epsilon `$-expansion using Bessel functions we consider a toy model integral related to the one-loop case in two dimensions. We select this example because the integral is finite and analytically known, so that we can compare our numerical calculation with the exact answer. Using Eq. (27), we obtain
$$2K_\epsilon (x)K_\epsilon (x)x๐x=\mathrm{\Gamma }(1+\epsilon )\mathrm{\Gamma }(1\epsilon ).$$
(71)
The expansion in $`\epsilon `$ is given by
$$\mathrm{\Gamma }(1\epsilon )\mathrm{\Gamma }(1+\epsilon )=\frac{\pi \epsilon }{\mathrm{sin}(\pi \epsilon )}=1+\frac{\pi ^2\epsilon ^2}{6}+\frac{7\pi ^4\epsilon ^4}{360}+O(\epsilon ^6)$$
(72)
On the other hand, we can use the expansion
$$K_\epsilon (x)=f_0(x)+\frac{\epsilon ^2}{2}f_2(x)+\frac{\epsilon ^4}{24}f_4(x)+O(\epsilon ^6)$$
(73)
to rewrite the integral in the form
$`2{\displaystyle _0^{\mathrm{}}}K_\epsilon (x)K_\epsilon (x)x๐x`$ $`=`$ $`2{\displaystyle _0^{\mathrm{}}}f_0(x)^2x๐x+2\epsilon ^2{\displaystyle _0^{\mathrm{}}}f_0(x)f_2(x)x๐x`$ (74)
$`+{\displaystyle \frac{\epsilon ^4}{6}}{\displaystyle _0^{\mathrm{}}}f_0(x)f_4(x)x๐x+{\displaystyle \frac{\epsilon ^4}{2}}{\displaystyle _0^{\mathrm{}}}f_2(x)^2x๐x+O(\epsilon ^6).`$
Using the explicit expressions for the functions $`f_k`$ we have checked by numerical integration that the identities
$`2{\displaystyle _0^{\mathrm{}}}f_0(x)^2x๐x`$ $`=`$ $`1,`$
$`2{\displaystyle _0^{\mathrm{}}}f_0(x)f_2(x)x๐x`$ $`=`$ $`{\displaystyle \frac{\pi ^2}{6}},`$
$`{\displaystyle \frac{1}{6}}{\displaystyle _0^{\mathrm{}}}(f_0(x)f_4(x)+3f_2(x)^2)x๐x`$ $`=`$ $`{\displaystyle \frac{7\pi ^4}{360}}`$ (75)
are valid numerically with very high degree of accuracy. We have implemented our algorithm for the $`\epsilon `$-expansion of sunrise-type diagrams as a simple code in Wolframโs Mathematica system for symbolic manipulations and checked its work-ability and efficiency.
When analyzing the $`\epsilon `$-expansion one realizes that integrals over products of modified Bessel functions of the second kind result in integrals over products of functions $`f_k(z)`$. But because the analytical behaviour of the functions $`f_k(z)`$ is quite similar to the original Bessel functions, the numerical integration is again easy to perform. For example, in the case of the four-loop bubble diagram, the integral
$$\stackrel{~}{\mathrm{\Pi }}_4(p^2=0)=D(x,m)^5d^Dx$$
(76)
can be calculated by subtracting the expansion (cf. Appendix E)
$`\mathrm{\Delta }(x)`$ $`=`$ $`{\displaystyle \frac{1}{4\pi ^{\lambda +1}x^{2\lambda }}}\{\mathrm{\Gamma }(\lambda )+(\left({\displaystyle \frac{x}{2}}\right)^2{\displaystyle \frac{\mathrm{\Gamma }(\lambda )}{1\lambda }}\left({\displaystyle \frac{x}{2}}\right)^{2\lambda }{\displaystyle \frac{\mathrm{\Gamma }(1\lambda )}{\lambda }})`$ (77)
$`+\left({\displaystyle \frac{x}{2}}\right)^2(\left({\displaystyle \frac{x}{2}}\right)^2{\displaystyle \frac{\mathrm{\Gamma }(\lambda )}{2(1\lambda )(2\lambda )}}{\displaystyle \frac{x}{2}}^{2\lambda }{\displaystyle \frac{\mathrm{\Gamma }(1\lambda )}{\lambda (\lambda +1)}})\}`$
from each of the propagators except for one which we keep as IR regulator. We obtain
$`\stackrel{~}{\mathrm{\Pi }}_4(0)`$ $`=`$ $`{\displaystyle D(x,m)\left(D(x,m)\mathrm{\Delta }(x)+\mathrm{\Delta }(x)\right)^4d^Dx}`$ (78)
$`=`$ $`{\displaystyle }D(x,m)((D(x,m)\mathrm{\Delta }(x))^4+4(D(x,m)\mathrm{\Delta }(x))^3\mathrm{\Delta }(x)`$
$`+6(\mathrm{\Delta }(x,m)\mathrm{\Delta }(x))^2\mathrm{\Delta }(x)^2+4(D(x,m)\mathrm{\Delta }(x))\mathrm{\Delta }(x)^3+\mathrm{\Delta }(x)^4)d^Dx.`$
The last two terms in Eq. (78) can be integrated analytically and can be expanded to an arbitrary order in $`\epsilon `$. They contain the singular contributions, i.e. all poles in $`\epsilon `$ and are expressible through Eulerโs $`\mathrm{\Gamma }`$-functions. As expected, the pole part coincides with the expression in Eq. (3.4). Because the analytical expression (as given up to order $`\epsilon ^3`$ in Appendix F) is rather lengthy, we present only its numerical evaluation,
$`๐ฉ_4\stackrel{~}{\mathrm{\Pi }}_4^{\mathrm{ana}}(0)`$ $`=`$ $`m^6(2.5\epsilon ^411.6666667\epsilon ^331.701389\epsilon ^267.528935\epsilon ^1`$ (79)
$`15871.965743142923.10240\epsilon 701868.64762\epsilon ^2`$
$`2486982.5547\epsilon ^3+O(\epsilon ^4)).`$
The remaining first three terms in Eq. (78) can be integrated numerically for $`D=4`$ using the expansion of the Bessel functions in terms of the functions $`a_k(z)`$ and $`b_k(z)`$ . No regularization is necessary since these contributions are regular at small $`x`$. The analytical expression for the functions to be integrated is again rather long. The zeroth order $`\epsilon `$-coefficient is found in Appendix G (see also the discussion of the integration procedure given in Appendix G). However, as shown in Fig. 5, the integrands themselves show a very smooth behaviour which renders the numerical integration rather simple. We obtain
$$๐ฉ_4\stackrel{~}{\mathrm{\Pi }}_4^{\mathrm{num}}(0)=m^6\left(15731.745122+142349.56687\epsilon +699112.42072\epsilon ^2+2468742.6339\epsilon ^3+O(\epsilon ^4)\right).$$
(80)
The sum of both parts gives
$`๐ฉ_4\stackrel{~}{\mathrm{\Pi }}_4(0)`$ $`=`$ $`m^6(2.5\epsilon ^411.6666667\epsilon ^331.701389\epsilon ^267.528935\epsilon ^1`$ (81)
$`140.220621573.53553\epsilon 2756.22690\epsilon ^218239.9208\epsilon ^3+O(\epsilon ^4))`$
which confirms the known result . We have not included as many digits as are given in Ref. .
It is not difficult to extend the analysis to higher orders in $`\epsilon `$ or to a larger number of significant digits in the coefficients of the $`\epsilon `$-expansion. However, since the technique is rather straightforward and simple we do not consider it worthwhile to extend the calculations into these directions. If the need arises, the potential user can tailor and optimize his or her programming code to obtain any desired accuracy and/or order of the $`\epsilon `$-expansion. In our evaluation we have used standard tools provided by Wolframโs Mathematica system for symbolic manipulations which allows one to reliably control the accuracy of numerical calculations. Even at this early stage of improvement it is obvious that our algorithm is extremely simple and efficient.
### 3.8 $`\epsilon `$-expansion of the five-loop vacuum bubble
In this paragraph we present results for the next-order sunrise-type diagram with $`p^2=0`$, the five-loop vacuum bubble (see Fig. 6). We have chosen to extend our calculation to the five-loop case since there exist no results on the five-loop bubble in the literature. The integral representation of the five-loop bubble with equal masses $`m`$ is given by
$$\stackrel{~}{\mathrm{\Pi }}_5(p^2=0)=D(x,m)^6d^Dx.$$
(82)
Evaluating numerically the analytical part one obtains
$`๐ฉ_5\stackrel{~}{\mathrm{\Pi }}_5^{\mathrm{ana}}(0)`$ $`=`$ $`m^8(3\epsilon ^5+16.5\epsilon ^4+51.95833\epsilon ^3+125.6715\epsilon ^2+259.9876\epsilon ^1`$ (83)
$`1360392.593416888723.177\epsilon 111392297.46\epsilon ^2`$
$`518606741.1\epsilon ^3+O(\epsilon ^4))`$
while the numerical integration of the nonsingular part gives
$`๐ฉ_5\stackrel{~}{\mathrm{\Pi }}_5^{\mathrm{num}}(0)`$ $`=`$ $`m^8(1360739.9485+16886269.683\epsilon `$ (84)
$`+111360751.91\epsilon ^2+518295438.0\epsilon ^3+O(\epsilon ^4)).`$
The sum of both contributions is given by
$`๐ฉ_5\stackrel{~}{\mathrm{\Pi }}_5(0)`$ $`=`$ $`m^8(3\epsilon ^5+16.5\epsilon ^4+51.95833\epsilon ^3+125.6715\epsilon ^2+259.9876\epsilon ^1`$ (85)
$`+347.35512453.494\epsilon 31545.55\epsilon ^2311303.1\epsilon ^3+O(\epsilon ^4)).`$
One observes huge cancellation effects between the terms obtained by the analytical calculation and the numerical integration. Apparently the subtraction procedure chosen here is non-optimal. As mentioned before, the subtraction procedure should really be optimized for any given problem in order to avoid a necessity to retain high numerical precision at intermediate steps of the calculation. Nevertheless, our non-optimized simple subtraction procedure already works quite reliably with available standard computational tools.
### 3.9 The spectral density, the discontinuity across the physical cut and the evaluation of the multi-particle phase space
Finally, we discuss the spectral density which can also be interpreted as the phase space volume of a multi-body decay of a single particle. An example of such a multi-body decay is shown in Fig. 7. At a given value of $`s`$ the value of the spectral density gives the phase space volume of the system of particles corresponding to the lines of the sunrise diagram. The phase space volume for $`(n+1)`$ particles with masses $`m_1,\mathrm{},m_{n+1}`$ is computed according to
$$\mathrm{\Phi }_{n+1}(p;p_1,\mathrm{},p_{n+1})=\frac{d^3p_1}{(2\pi )^3E_1}\mathrm{}\frac{d^3p_{n+1}}{(2\pi )^32E_{n+1}}\delta ^{(4)}(p_1+\mathrm{}+p_{n+1}p)$$
where $`E_i=\sqrt{\stackrel{}{p}_i^2+m_i^2}`$. The optical theorem relates the imaginary part of a Feynman amplitude to on-shell propagation of its internal lines and therefore to the product of an outer state with its conjungate via so-called cutting rules (for a recent review see also Ref. ). Therefore the discontinuity across a physical cut of the sunrise-type diagram with $`(n+1)`$ internal lines carrying $`(n+1)`$ masses $`m_1,\mathrm{},m_{n+1}`$ leads to the $`(n+1)`$ particle phase space in a natural way. We obtain
$$\mathrm{\Phi }_{n+1}(p;p_1,\mathrm{},p_{n+1})=2\pi \rho (s)|_{p^2=s}$$
which means that, up to a constant factor, the phase space is given by the spectral density.
## 4 Asymptotic analysis of sunrise-type diagrams
In multi-loop calculations it is very useful to use expansions in different regimes of its kinematical variables and masses, corresponding each to the smallness of one of the parameters involved in the problem. After defining an appropriate scale, the separation of scales can considerably simplify the description in the given region of energy or momentum transfer for a particular physical process. Such an analysis is needed when, starting from an underlying full theory, one wants to formulate an effective theory for a given limited energy scale. Taking for instance Quantum Chromodynamics (QCD) as the underlying theory, the best-known of such effective theories is the Heavy Quark Effective Theory (HQET) which analyzes the behaviour of a heavy quark near its mass shell which leads to a simplified description of physical processes subject to this approximation . The effective theory is directly constructed through an expansion of the expression originally written in QCD. Its analysis, therefore, relies on the expansion of Feynman diagrams in a special regime. Another important example is the expansion near the production threshold, leading to an effective theory termed nonrelativistic QCD (NRQCD). NRQCD is a more complicated case of an effective theory since one introduces new variables in NRQCD which differ from those defined in QCD. The most involved example is Chiral Perturbation Theory (ChPT) where the variables are Goldstone modes which are not directly related to QCD (and where it is still difficult to construct a direct correspondence). Nevertheless, the common thread of all these effective theories is that they use perturbative expansions in terms of diagrams in given specific kinematical regimes.
When sunrise-type diagrams enter these analyses, they usually provide the leading order contribution or serve as laboratory for checking more elaborate techniques. In this connection the expansion of sunrise-type diagrams in configuration space happens to be very efficient for many different regimes. Having the closed form expression in Eq. (18) for the correlation function $`\stackrel{~}{\mathrm{\Pi }}(p)`$ at hand, different asymptotic regimes in the parameters (masses, external momentum) can easily be analyzed by calculating the corresponding expansions. The first three subsections in this section will deal with examples in the Euclidean domain. Starting with the fourth subsection, we will consider some limiting cases in the Minkowskian domain which is the domain of physical states. The calculation of asymptotic cases is more involved in this case. The calculations are performed for the spectrum given by Eq. (24) which contains the physical content of the problem.
### 4.1 Large momentum expansion: close to the massless limit
Loop integrations considerably simplify for massless internal lines. This holds true for the momentum space representation as well as for the configuration space representation of the loops. In some cases particle mass corrections to the massless approximation can become numerically important. For instance, in QCD the analysis of hadronic processes involving the strange quark require some special care since the mass of the $`s`$-quark is not small compared to the masses of the lightest quarks. On the other hand, the $`s`$-quark mass is still much smaller than the QCD scale $`\mathrm{\Lambda }_{\mathrm{QCD}}`$. The expansion in the ratio of the $`s`$-quark mass to the relevant scale of the considered process is therefore very effective. This expansion is especially fruitful for the analysis of tau lepton decays .
Mass corrections to the large $`p^2`$ behaviour in the Euclidean domain (i.e. expansions in $`m_i^2/p^2`$) are obtained by expanding the massive propagators under the integration sign in Eq. (18) in terms of the masses $`m_i`$. The final integration is performed by using the identity
$$_0^{\mathrm{}}x^\mu J_\lambda (px)๐x=2^\mu p^{\mu 1}\frac{\mathrm{\Gamma }\left((\lambda +\mu +1)/2\right)}{\mathrm{\Gamma }\left((\lambda \mu +1)/2\right)}.$$
(86)
Note that all these manipulations are straightforward and can be easily implemented in a system of symbolic computations. Some care is necessary, however, when poles of the $`\mathrm{\Gamma }`$-function are encountered which reflect the presence of artificial infrared singularities. The framework for dealing with such problems is well-known (see e.g. Refs. ).
A further physics application concerns the study of pentaquark properties within the QCD sum rule approach. In Fig. 8 we show the leading order contribution to the relevant two-point correlator of the interpolating pentaquark currents which contains one massive $`s`$-quark line. One therefore has to determine the mass corrections to the two-point correlator. Another application of the proposed techniques concerns the analysis of certain scalar mesons which can be viewed as four-quark states involving the $`s`$-quark. For instance, the scalar meson $`a_0`$ can be considered to be a $`KK`$ molecule built from two $`K`$-mesons . Yet another example is the analysis of $`K^0\overline{K}^0`$ mixing in the two-point function approach . Note, however, that the approach based on the three-point function analysis appears to be more direct .
### 4.2 Small momentum expansion: close to the vacuum bubble
The basic expression in Eq. (18) is also well suited for finding the power series expansion in $`p^2`$ because the values of the polarization function and its derivatives at $`p^2=0`$ can easily be obtained. The convenience of a $`p^2`$-expansion is demonstrated by making use of the basic identity for differentiating the Bessel function (cf. Eq. (A6)),
$$\frac{d^k}{d(p^2)^k}\left(\frac{px}{2}\right)^\lambda J_\lambda (px)=\left(\frac{x^2}{4}\right)^k\left(\frac{px}{2}\right)^{\lambda k}J_{\lambda +k}(px).$$
(87)
Differentiation, therefore, results in an expression which has the same functional structure as the original function. This is convenient for numerical computations. Note that sufficiently high order derivatives become UV-finite. Operationally, this is obvious since the subtraction polynomial vanishes after sufficiently high derivatives are taken. However, this can also be seen explicitly from Eq. (87) where high powers of $`x^2`$ suppress the singularity of the product of propagators at small $`x`$. The key integral is given by
$$_0^{\mathrm{}}x^{\mu 1}K_\nu (Mx)๐x=2^{\mu 2}M^\mu \mathrm{\Gamma }\left(\frac{\mu +\nu }{2}\right)\mathrm{\Gamma }\left(\frac{\mu \nu }{2}\right).$$
(88)
The main physical applications are given by calculations within ChPT and HQET. The technique is especially useful for the analysis of form factors. Another application is given by expansions of heavy quark correlators through the conformal mapping in order to find an approximate spectral density from the calculation of moments . Finally, applications can be found in the analysis of some specific rules for massive exotic states by calculating the moments of the spectral density.
### 4.3 Dominating mass expansion: near the static limit
Consider a hadron with one large mass which is of the same order as the external momentum and all other masses are small. An example is a heavy meson consisting of the $`c`$-quark and light quark(s). In the Euclidean domain the polarization function is not difficult to compute in this limit because of the simplicity of the configuration space representation and the high speed of convergence of the ensuing numerical procedures. One can do even better since this special limit can be done analytically. When expanding the propagators in the limit of small masses one encounters powers of $`x`$ and $`\mathrm{ln}(mx)`$ (or, within the framework of dimensional regularization, non-integer powers of $`x`$). The remaining functions are the weight function (Bessel function of the first kind) and the propagator of the heavy particle with the large mass $`M`$ which is given by the modified Bessel function of the second kind. The general structure of the terms in the series expansion that contribute to $`\stackrel{~}{\mathrm{\Pi }}(p)`$ is given by
$$2\pi ^{\lambda +1}_0^{\mathrm{}}\left(\frac{px}{2}\right)^\lambda J_\lambda (px)K_\nu (Mx)x^{2\rho }๐x$$
(89)
The integrations in Eq. (89) can be done in closed form by using the basic integral representation
$`{\displaystyle _0^{\mathrm{}}}x^\mu J_\lambda (px)K_\nu (Mx)๐x`$ (90)
$`=`$ $`{\displaystyle \frac{p^\lambda \mathrm{\Gamma }((\lambda +\mu +\nu +1)/2)\mathrm{\Gamma }((\lambda +\mu \nu +1)/2)}{2^{1\mu }M^{\lambda +\mu +1}\mathrm{\Gamma }(\lambda +1)}}`$
$`\times _2F_1((\lambda +\mu +\nu +1)/2,(\lambda +\mu \nu +1)/2;\lambda +1;p^2/M^2)`$
where $`{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ is the hypergeometric function . The corresponding integrals contain integer powers of logarithms. They can be computed by differentiation with respect to $`\mu `$. Note that the maximal power of the logarithm is determined by the number of light propagators and does not increase with the order of the expansion: any light propagator contains only one power of the logarithm, as can be seen from the expansion of the Bessel function $`K_\nu (\xi )`$ at small $`\xi `$ .
Physical applications in this regime are mainly sum rules for hadrons including a heavy flavour . These are the sum rules for the $`\mathrm{\Lambda }_c`$ or $`\mathrm{\Lambda }_b`$ baryons where the lowest order contribution to the sum rules are given by the genuine two-loop sunrise-type diagram. In the case of the $`D`$ and $`B`$ mesons only the one-loop (degenerate sunrise) diagram contributes.
### 4.4 Threshold expansion: close to nonrelativistic physics
The threshold region is very important for the description of heavy quarkonia states. In the threshold region the necessary calculations considerably simplify if one uses an appropriate effective theory. This effective theory is designed to remove the degrees of freedom that can be treated perturbatively from the outset and considers only the dynamics of the essential modes relevant for the considered energy scales. For heavy quarkonia near the threshold the effective theory is constructed on the basis of the non-relativistic approximation for the strong interaction and is called nonrelativistic QCD (NRQCD). Due to the use of NRQCD there has been a great advance in describing the heavy flavour production near threshold during the last years .
The threshold region of a sunrise-type diagram is determined by the condition $`q^2+M^20`$ where $`q`$ is the Euclidean momentum and $`M=_im_i`$ is the threshold value for the spectral density. We introduce the Minkowskian momentum $`p`$ defined by $`p^2=q^2`$ which allows for an analytic continuation to the physical cut. Operationally this analytic continuation can be performed by replacing $`qip`$. To analyze the region near threshold we use the parameter $`\mathrm{\Delta }=Mp`$ which can take complex values. The parameter $`\mathrm{\Delta }`$ is more convenient in the Euclidean domain while the parameter $`E=\mathrm{\Delta }=pM`$ is the actual energy counted from threshold and is used in phenomenological applications. The analytic continuation of the Fourier transform to the Minkowskian domain has the form
$$\stackrel{~}{\mathrm{\Pi }}(p)=2\pi ^{\lambda +1}_0^{\mathrm{}}\left(\frac{ipx}{2}\right)^\lambda J_\lambda (ipx)\mathrm{\Pi }(x)x^{2\lambda +1}๐x.$$
(91)
For the threshold expansion we have to analyze the large $`x`$ behaviour of the integrand, corresponding to the region which saturates the integral in the limit $`pM`$ or, equivalently, $`E0`$. It is convenient to perform the analysis in a basis where the integrand has a simple large $`x`$ behaviour. The most important part of the integrand is the Bessel function $`J_\lambda (ipx)`$ which, however, contains both rising and falling branches at large $`x`$. This resembles the situation with the elementary trigonometric functions $`\mathrm{sin}(z)`$ and $`\mathrm{cos}(z)`$ to which the Bessel function $`J_\lambda (z)`$ is rather close in a certain sense. Indeed, the function $`\mathrm{cos}(z)`$ and $`\mathrm{sin}(z)`$ are linear combinations of exponentials, namely
$$\mathrm{cos}(z)=\frac{1}{2}\left(e^{iz}+e^{iz}\right)$$
(92)
and has also both rising and falling branches at large pure imaginary arguments: the exponentials show simple asymptotic behaviour $`e^{\pm z}`$ at $`z=\pm i\mathrm{}`$. The analogous statement is true for $`J_\lambda (z)`$ which can be written as a sum of two Hankel functions,
$$J_\lambda (z)=\frac{1}{2}(H_\lambda ^+(z)+H_\lambda ^{}(z))$$
(93)
where $`H_\lambda ^\pm (z)=J_\lambda (z)\pm iY_\lambda (z)`$. The Hankel functions $`H_\lambda ^\pm (z)`$ show a simple asymptotic behaviour at infinity,
$$H_\lambda ^\pm (iz)z^{1/2}e^{\pm z}$$
(94)
(cf. Eq. (A16)). Accordingly we split up $`\stackrel{~}{\mathrm{\Pi }}(p)`$ into $`\stackrel{~}{\mathrm{\Pi }}(p)=\stackrel{~}{\mathrm{\Pi }}^+(p)+\stackrel{~}{\mathrm{\Pi }}^{}(p)`$ with
$$\stackrel{~}{\mathrm{\Pi }}^\pm (p)=\pi ^{\lambda +1}_0^{\mathrm{}}\left(\frac{ipx}{2}\right)^\lambda H_\lambda ^\pm (ipx)\mathrm{\Pi }(x)x^{2\lambda +1}๐x.$$
(95)
The two parts $`\stackrel{~}{\mathrm{\Pi }}^\pm (p)`$ of the polarization function $`\stackrel{~}{\mathrm{\Pi }}(p)`$ have a completely different behaviour near threshold which allows one to analyze them independently. This observation makes the subsequent analysis straightforward.
We first consider the contribution of the $`\stackrel{~}{\mathrm{\Pi }}^+(p)`$ part which reduces to a regular sunrise-type diagram. Indeed, using the relation given by Eq. (A14),
$$K_\lambda (z)=\frac{\pi i}{2}e^{i\lambda \pi /2}H_\lambda ^+(iz)$$
(96)
between Bessel functions of different kinds one can replace the Hankel function $`H_\lambda ^+(ipx)`$ with the Bessel function $`K_\lambda (px)`$. Since the propagator of a massive particle (massive line in the diagram) is given by the Bessel function up to a power in $`x`$, this substitution shows that the weight function behaves like a propagator of an additional line with the โmassโ $`p`$. The explicit expression is given by
$$\stackrel{~}{\mathrm{\Pi }}^+(p)=\frac{(2\pi i)^{2\lambda +1}}{(p^2)^\lambda }_0^{\mathrm{}}\mathrm{\Pi }_+(x)x^{2\lambda +1}๐x.$$
(97)
The function $`\mathrm{\Pi }_+(x)=\mathrm{\Pi }(x)D(x,p)`$ is the polarization function of a new effective diagram which is equal to the initial polarization function multiplied by a propagator with $`p`$ as mass parameter. We thus end up with a vacuum bubble of the sunrise-type type with one additional line compared to the initial diagram (see Fig. 9). All derivatives of $`\stackrel{~}{\mathrm{\Pi }}^+(p)\stackrel{~}{\mathrm{\Pi }}^+(M\mathrm{\Delta })`$ with respect to $`\mathrm{\Delta }`$ are represented as vacuum bubbles with one additional line carrying rising indices. Such diagrams can be efficiently calculated within the recurrence relation technique developed in . They possess no singularities at the production threshold $`p=M`$. This can be seen by looking at the expansion for large $`x`$. The behaviour at large $`x`$ is given by the asymptotic form of the functions for which one obtains
$$H^+(ipx)=\sqrt{\frac{2}{i\pi px}}e^{px}(1+O(x^1)),K(mx)=\sqrt{\frac{\pi }{2mx}}e^{mx}(1+O(x^1))$$
(98)
(cf. Eqs. (A15) and (A16)). The large $`x`$ range of the integral (above a reasonably large cutoff parameter $`\mathrm{\Lambda }`$) has the general form
$$\stackrel{~}{\mathrm{\Pi }}_\mathrm{\Lambda }^+(M\mathrm{\Delta })_\mathrm{\Lambda }^{\mathrm{}}x^ae^{(2M\mathrm{\Delta })x}๐x$$
(99)
where
$$a=(n1)(\lambda +1/2).$$
(100)
The right hand side of Eq. (99) is an analytic function in $`\mathrm{\Delta }`$ in the vicinity of $`\mathrm{\Delta }=0`$. It exhibits no cut or other singularities near the threshold and therefore does not contribute to the spectral density.
In contrast to the previous case, the integrand of the second part $`\stackrel{~}{\mathrm{\Pi }}^{}(p)`$ contains $`H^{}(ipx)`$ which behaves like a rising exponential function at large $`x`$,
$$H^{}(ipx)x^{1/2}e^{px}$$
(101)
(cf. Eq. (A16). Therefore, the integral is represented by
$$\stackrel{~}{\mathrm{\Pi }}_\mathrm{\Lambda }^{}(M\mathrm{\Delta })_\mathrm{\Lambda }^{\mathrm{}}x^ae^{\mathrm{\Delta }x}๐x.$$
(102)
The function $`\stackrel{~}{\mathrm{\Pi }}^{}(M\mathrm{\Delta })`$ is non-analytic near $`\mathrm{\Delta }=0`$ because for $`\mathrm{\Delta }<0`$ the integrand in Eq. (102) grows in the large $`x`$ region and the integral diverges at the upper limit. Therefore the function which is determined by this integral is singular at $`\mathrm{\Delta }<0`$ ($`E>0`$) and requires an interpretation for these values of the argument $`\mathrm{\Delta }`$. The function is analytic in the complex $`\mathrm{\Delta }`$-plane with a cut along the negative axis. This cut corresponds to the physical positive energy cut. The discontinuity across the cut gives rise to the non-vanishing spectral density of the diagram (cf. Appendix C).
The integral for $`\stackrel{~}{\mathrm{\Pi }}^{}(p)`$ in Eq. (95) cannot be done analytically. In order to obtain an expansion for the spectral density near the threshold in analytical form we make use of the asymptotic series expansion for the function $`\mathrm{\Pi }(x)`$ which crucially simplifies the integrands but still preserves the singular structure of the integral in terms of the variable $`\mathrm{\Delta }`$. The asymptotic series expansion of the main part of each propagator, i.e. of the modified Bessel function of the second kind, to the order $`N`$ is given by Eq. (A15). The asymptotic expansion of the function $`\mathrm{\Pi }(x)`$ therefore consists of an exponential factor $`e^{Mx}`$ and an inverse power series in $`x`$ up to the order $`\stackrel{~}{N}`$, where $`\stackrel{~}{N}`$ is closely related to $`N`$. It is this asymptotic expansion that determines the singularity structure of the integral. We write the whole integral in the form of a sum of two terms,
$`\stackrel{~}{\mathrm{\Pi }}^{}(p)`$ $`=`$ $`\pi ^{\lambda +1}{\displaystyle \left(\frac{ipx}{2}\right)^\lambda H_\lambda ^{}(ipx)\left(\mathrm{\Pi }(x)\mathrm{\Pi }_N^{as}(x)\right)x^{2\lambda +1+2\epsilon }๐x}`$
$`+\pi ^{\lambda +1}{\displaystyle \left(\frac{ipx}{2}\right)^\lambda H_\lambda ^{}(ipx)\mathrm{\Pi }_N^{as}(x)x^{2\lambda +1+2\epsilon }๐x}=\stackrel{~}{\mathrm{\Pi }}^{di}(p)+\stackrel{~}{\mathrm{\Pi }}^{as}(p).`$
The integrand of the first term $`\stackrel{~}{\mathrm{\Pi }}^{di}(p)`$ behaves as $`1/x^{\stackrel{~}{N}}`$ at large $`x`$ while the integrand of the second term accumulates all lower powers of the large $`x`$ expansion. Note that only the large $`x`$ behaviour is essential for the near threshold expansion of the spectral density. This fact has been taken into account in Eqs. (99) and (102) where we introduced a cutoff $`\mathrm{\Lambda }`$. However, from the practical point of view the calculation of the regularized integrals with an explicit cutoff is inconvenient. The final result of the calculation โ the spectral density of the diagram โ is independent of the cutoff, but the integration is technically complicated if the cutoff is introduced. However, in extending the integration over the whole region of the variable $`x`$ without using the cutoff one immediately encounters divergences at small $`x`$ because the asymptotic expansion is invalid in the region near the origin, so one is not allowed to continue it to this region. The standard way to cope with such a situation is to introduce dimensional regularization. It allows one to deal with divergent expressions at intermediate stages of the calculation and is technically simple because it does not introduce any cutoff and therefore does not modify the integration region very much. Note that dimensional regularization does not necessarily regularize all divergences in this case (in contrast to the standard case of ultraviolet divergences) but it nevertheless suffices for our purposes. We shall use a parameter $`\epsilon `$ to regularize the divergences at small $`x`$.
The first part $`\stackrel{~}{\mathrm{\Pi }}^{di}(p)`$ in Eq. (4.4) containing a difference of the polarization function and its asymptotic expansion since the integrand gives no contributions to the spectral density up to a given order of the expansion in $`\mathrm{\Delta }`$. This is so because the subtracted asymptotic series to order $`N`$ cancels the inverse power behaviour of the integrand to this degree $`N`$. The integrand decreases sufficiently fast for large values of $`x`$ and the integral converges even at $`\mathrm{\Delta }=0`$. Therefore this term is inessential for the evaluation of the expansion of the spectral density up to given order.
The expansion of the spectral density at small $`E`$ is determined only by the integral $`\stackrel{~}{\mathrm{\Pi }}^{as}(p)`$ in Eq. (4.4). This integral is still rather complicated to compute but we can go a step further in its analytical evaluation. Indeed, since the singular behaviour of $`\stackrel{~}{\mathrm{\Pi }}^{as}(p)`$ is determined by the behaviour at large $`x`$, we can replace the first factor, i.e. the Hankel function, in the large $`x`$ region by its asymptotic expansion up to some order $`N`$. We use Eq. (A16) to obtain a representation
$$\stackrel{~}{\mathrm{\Pi }}^{das}(p)=\pi ^{\lambda +1}\left(\frac{ipx}{2}\right)^\lambda H_{\lambda ,N}^{as}(ipx)\mathrm{\Pi }_N^{as}(x)x^{2\lambda +1+2\epsilon }๐x.$$
(104)
The index โdasโ stands for โdouble asymptoticโ and indicates that the integrand in Eq. (104) consists of a product of two asymptotic expansions: one for the polarization function $`\mathrm{\Pi }(x)`$ and another for the Hankel function $`H_\lambda (x)`$ as weight (or kernel). Both asymptotic expansions are straightforward and can be obtained from standard handbooks on Bessel functions (cf. Appendix A). We therefore arrive at our final result: the integration necessary for evaluating the near threshold expansion of the sunrise-type diagrams reduces to integrals of Eulerโs Gamma function type, i.e. integrals containing exponentials and powers. Indeed, the result of the expansion in Eq. (104) is an exponential function $`e^{\mathrm{\Delta }x}`$ times a power series in $`1/x`$, namely
$$x^{a+2\epsilon }e^{\mathrm{\Delta }x}\underset{j=0}{\overset{N1}{}}\frac{A_j}{x^j}$$
(105)
where $`a`$ has already been defined in Eq. (100) and the coefficients $`A_j`$ are simple functions of the momentum $`p`$ and the masses $`m_i`$. The expression in Eq. (105) can be integrated analytically using
$$_0^{\mathrm{}}x^{a+2\epsilon }e^{\mathrm{\Delta }x}๐x=\mathrm{\Gamma }(1a+2\epsilon )\mathrm{\Delta }^{a12\epsilon }.$$
(106)
The result is
$$\stackrel{~}{\mathrm{\Pi }}^{das}(M\mathrm{\Delta })=\underset{j=0}{\overset{N1}{}}A_j\mathrm{\Gamma }(1aj+2\epsilon )\mathrm{\Delta }^{a+j12\epsilon }.$$
(107)
This expression is our final representation for the part of the polarization function of a sunrise-type diagram necessary for the calculation of the spectral density near the production threshold (see Appendix C). The spectral density is a function of $`E=\mathrm{\Delta }`$ and will be denoted by $`\stackrel{~}{\rho }(E)=\rho ((M+E)^2)`$ in the following. Starting from the main result in Eq. (107), we discuss the general structure in detail. In the case where $`a`$ takes integer values, these coefficients result in $`1/\epsilon `$-divergences for small values of $`\epsilon `$. The powers of $`\mathrm{\Delta }`$ in Eq. (107) have to be expanded to first order in $`\epsilon `$ and give
$$\frac{1}{2\epsilon }\mathrm{\Delta }^{2\epsilon }=\frac{1}{2\epsilon }+\mathrm{ln}\mathrm{\Delta }+O(\epsilon ).$$
(108)
Because of
$$\mathrm{Disc}\mathrm{ln}(\mathrm{\Delta })\mathrm{ln}(Ei0)\mathrm{ln}(E+i0)=2\pi i\theta (E),$$
(109)
$`\stackrel{~}{\mathrm{\Pi }}^{das}(M\mathrm{\Delta })`$ in Eq. (107) contributes to the spectral density. For half-integer values of $`a`$ the power of $`\mathrm{\Delta }`$ itself has a cut even for $`\epsilon =0`$. The discontinuity is then given by
$$\mathrm{Disc}\sqrt{\mathrm{\Delta }}=2i\sqrt{E}\theta (E).$$
(110)
Our method to construct a threshold expansion thus simply reduces to the analytical calculation of the integral in Eq. (104) which can be done for arbitrary dimension and an arbitrary number of lines with different masses. In the next paragraphs we use our technique to work out some specific examples which demonstrate both the simplicity and efficiency of our method.
### 4.5 Equal mass genuine sunrise diagram at threshold
The polarization function corresponding to the genuine sunrise diagram with three propagators with equal masses $`m`$ in $`D=4`$ space-time dimensions is given by
$$\mathrm{\Pi }(x)=\frac{m^3K_1(mx)^3}{(2\pi )^6x^3}.$$
(111)
The exact spectral density is given by the integral representation in Eq. (24) which for this particular case reads
$$\rho (s)=\frac{2\pi }{i\sqrt{s}}_{ci\mathrm{}}^{c+i\mathrm{}}I_1(x\sqrt{s})\mathrm{\Pi }(x)x^2๐x.$$
(112)
In order to obtain a threshold expansion of the spectral density in Eq. (112) we use Eq. (107) to calculate the expansion of the appropriate part of the polarization function. To illustrate the procedure we derive the explicit functional behaviour of the integrand in Eq. (104) which is given by an asymptotic expansion at large $`x`$,
$`\pi ^2\left({\displaystyle \frac{ipx}{2}}\right)^1H_{1,N}^{as}(px)\mathrm{\Pi }_N^{as}(x)x^{3+2\epsilon }={\displaystyle \frac{m^{3/2}e^{(p3m)x}}{(4\pi )^3p^{3/2}}}x^{3+2\epsilon }\times `$ (113)
$`\times \left\{1+{\displaystyle \frac{9}{8mx}}{\displaystyle \frac{3}{8px}}+{\displaystyle \frac{9}{128m^2x^2}}{\displaystyle \frac{27}{64mpx^2}}{\displaystyle \frac{15}{128p^2x^2}}+O(x^3)\right\}.`$
From Eq. (113) we can easily read off the coefficients $`A_j`$ that enter the expansion in Eq. (105). The spectral density is obtained by performing a term-by-term integration of the series in Eq. (113) and by evaluating the discontinuity across the cut along the positive energy axis $`E>0`$. The result reads
$`\stackrel{~}{\rho }(E)={\displaystyle \frac{E^2}{384\pi ^3\sqrt{3}}}\{1{\displaystyle \frac{1}{2}}\eta +{\displaystyle \frac{7}{16}}\eta ^2{\displaystyle \frac{3}{8}}\eta ^3+{\displaystyle \frac{39}{128}}\eta ^4{\displaystyle \frac{57}{256}}\eta ^5`$
$`+{\displaystyle \frac{129}{1024}}\eta ^6{\displaystyle \frac{3}{256}}\eta ^7{\displaystyle \frac{4047}{32768}}\eta ^8+{\displaystyle \frac{18603}{65536}}\eta ^9{\displaystyle \frac{248829}{524288}}\eta ^{10}+O(\eta ^{11})\}`$
where the notation $`\eta =E/M`$, $`M=3m`$ is used. The simplicity of the derivation is striking. By no cost it can be generalized to any number of lines, arbitrary masses, and any space-time dimension. The genuine equal mass sunrise has been chosen for definiteness only. It also allows us to compare our results with results available in the literature. Eq. (4.5) reproduces the expansion coefficients $`\stackrel{~}{a}_j`$ obtained in Ref. (the fourth column in Table 1 of Ref. ) by a direct integration in momentum space within the technique of region separation . In Fig. 10 the exact solution is shown together with expansions to various orders.
### 4.6 Three loops and a route to any number of loops at threshold
The sunrise-type diagrams with four or more propagators cannot be easily done by using the momentum space technique because it requires the multi-loop integration of entangled momenta. As emphasized before the configuration space technique allows one to immediately generalize our previous results to any number of lines (or loops) with no additional effort. Consider first the three-loop case of sunrise-type diagrams. The polarization function of the equal mass sunrise-type diagram with four propagators in $`D=4`$ space-time is given by
$$\mathrm{\Pi }(x)=\frac{m^4K_1(mx)^4}{(2\pi )^8x^4}.$$
(115)
The exact spectral density of this diagram can be obtained from Eq. (112) while the near threshold expansion can be found using Eq. (107). The expansion of the spectral density near threshold reads
$`\stackrel{~}{\rho }(E)={\displaystyle \frac{E^{7/2}M^{1/2}}{26880\pi ^5\sqrt{2}}}\{1{\displaystyle \frac{1}{4}}\eta +{\displaystyle \frac{81}{352}}\eta ^2{\displaystyle \frac{2811}{18304}}\eta ^3+{\displaystyle \frac{17581}{292864}}\eta ^4`$
$`+{\displaystyle \frac{1085791}{19914752}}\eta ^5{\displaystyle \frac{597243189}{3027042304}}\eta ^6+{\displaystyle \frac{4581732455}{12108169216}}\eta ^7{\displaystyle \frac{496039631453}{810146594816}}\eta ^8+O(\eta ^9)\}`$
where $`\eta =E/M`$ and $`M=4m`$ is the threshold value. One sees the difference with the previous three-line case. In Eq. (4.6) the cut corresponds to a square root branch cut while in the three-line case one has a logarithmic cut. One can easily figure out the reason for this by looking at the asymptotic structure of the integrand. For an even number of lines (i.e. odd number of loops) one has a square root branch cut, while for an odd number of lines (even number of loops) one has a logarithmic branch cut. This statement generalizes from $`D=4`$ to any even space-time dimension. In the general case the structure of the cut depends on the dimensionality of the space-time. The general formula for the $`n`$-loop case reads
$$\stackrel{~}{\rho }(E)E^{(\lambda +1/2)n1}(1+O(E)).$$
(117)
For $`D=4`$ space-time dimension (i.e. $`\lambda =1`$) we can verify the result of Ref. ,
$$\stackrel{~}{\rho }(E)E^{(3n2)/2}(1+O(E)).$$
(118)
Returning to Eq. (4.6) one has numerically
$`\stackrel{~}{\rho }(E)=8.596210^5E^{7/2}M^{1/2}\{1.0000.250\eta +0.230\eta ^2`$
$`0.154\eta ^3+0.060\eta ^4+0.055\eta ^50.197\eta ^6+0.378\eta ^70.612\eta ^8+O(\eta ^9)\}`$
where we have written down the coefficients up to three decimal places. It is difficult to say anything definite about the convergence of this series. By construction it is an asymptotic series. However, we stress that the practical (or explicit) convergence can always be checked by comparing series expansions like the one shown in Eq. (4.6) with the exact spectral density given in Eq. (112) by numerical integration.
We conclude that the spectral density of the sunrise-type diagram can be efficiently calculated within the configuration space technique. It does not matter whether one is aiming for the exact result or an expansion the configuration space technique which can readily deliver the desired result. The exact formula in Eq. (112) as well as the threshold expansion obtained from it can be used to calculate the spectral density for an arbitrarily large number of internal lines. We stress that the case of different masses does not lead to any complications within the configuration space technique: the exact formula in Eq. (24) and/or the near threshold expansion work equally well for any configuration of masses. We do not present plots for general cases of different masses because they are not very illuminating as one can only see the common threshold. However, there is an interesting kinematic regime for the different mass case which is important for certain applications which, to the best of our knowledge, has not been treated in the literature before. An analytical solution for the expansion of the spectral density in this regime is given in the next subsection.
### 4.7 New features of the threshold expansion: resummation of small mass effects for strongly asymmetric mass configurations
The main question in constructing expansions is their practical usefulness and the region of validity. The threshold expansion for equal or almost equal masses breaks down for instance for $`EM=m_i`$. However, if the masses are not equal, the region of the breakdown of the expansion is determined by the mass with the smallest numerical value. The simplest example where one can see this phenomenon is the analytical expression for the spectral density of the simple loop (degenerate sunrise-type diagram) with two different masses $`m_1`$ and $`m_2`$. In $`D=4`$ space-time dimensions (see e.g. Ref. ) one has
$$\stackrel{~}{\rho }(E)=\frac{\sqrt{E(E+2m_1)(E+2m_2)(E+2M)}}{(4\pi (M+E))^2}$$
(120)
where $`M=m_1+m_2`$. The threshold expansion is obtained by expanding the right hand side of Eq. (120) in $`E`$ for small values of $`E`$. If $`m_2`$ is much smaller than $`m_1`$, the expansion breaks down at $`E2m_2`$. The break-down of the series expansion can also be observed in more general cases. If one of the masses (which we call $`m_0`$) is much smaller than the other masses, the threshold expansion is only valid in a very limited region $`E<<2m_0`$.
To generalize the expansion and extend it to the region of $`EM`$ one has to treat the smallest mass exactly. In this case one can use a method which we call the resummation of small mass effects . This method will be explained in the following. We start with the representation
$$\stackrel{~}{\mathrm{\Pi }}^{pas}(p)=\pi ^{\lambda +1}\left(\frac{ipx}{2}\right)^\lambda H_{\lambda ,N}^{as}(ipx)\mathrm{\Pi }_{m_0}^{as}(x)x^{2\lambda +1+2\epsilon }๐x$$
(121)
which is the part of the polarization function contributing to the spectral density. The integrand in Eq. (121) has the form
$$\mathrm{\Pi }_{m_0}^{as}(x)=\mathrm{\Pi }_{n1}^{as}(x)D(m_0,x)$$
(122)
where the asymptotic expansions are substituted for all the propagators except for the one with the small mass $`m_0`$. This is indicated by the index โpasโ in Eq. (121) which stands for โpartial asymptoticโ. The main technical observation leading to the generalization of the expansion method is that $`\stackrel{~}{\mathrm{\Pi }}^{pas}(p)`$ is still analytically computable in a closed form. Indeed, the genuine integral to compute has the form
$`{\displaystyle _0^{\mathrm{}}}x^{\mu 1}e^{\stackrel{~}{\alpha }x}K_\nu (\beta x)๐x=`$ (123)
$`=`$ $`{\displaystyle \frac{\sqrt{\pi }(2\beta )^\nu }{(2\stackrel{~}{\alpha })^{\mu +\nu }}}{\displaystyle \frac{\mathrm{\Gamma }(\mu +\nu )\mathrm{\Gamma }(\mu \nu )}{\mathrm{\Gamma }(\mu +1/2)}}{}_{2}{}^{}F_{1}^{}({\displaystyle \frac{\mu +\nu }{2}},{\displaystyle \frac{\mu +\nu +1}{2}};\mu +{\displaystyle \frac{1}{2}};1{\displaystyle \frac{\beta ^2}{\stackrel{~}{\alpha }^2}})`$
where $`\stackrel{~}{\alpha }=\mathrm{\Delta }m_0`$ and $`\beta =m_0`$. The integral $`\stackrel{~}{\mathrm{\Pi }}^{pas}(p)`$ in Eq. (121) is thus expressible in terms of hypergeometric functions . For the construction of the spectral density one has to find the discontinuity of the right hand side of Eq. (123). There are several ways to do this. We proceed by applying the discontinuity operation to the integrand of the integral representation of the hypergeometric function. The resulting integrals are calculated again in terms of hypergeometric functions. Indeed,
$`{\displaystyle \frac{1}{2\pi i}}\mathrm{Disc}{\displaystyle _0^{\mathrm{}}}x^{\mu 1}e^{\alpha x}K_\nu (\beta x)๐x=`$ (124)
$`=`$ $`{\displaystyle \frac{2^\mu (\alpha ^2\beta ^2)^{1/2\mu }}{\alpha ^{1/2\nu }\beta ^\nu }}{\displaystyle \frac{\mathrm{\Gamma }(3/2)}{\mathrm{\Gamma }(3/2\mu )}}{}_{2}{}^{}F_{1}^{}({\displaystyle \frac{1\mu \nu }{2}},{\displaystyle \frac{2\mu \nu }{2}};{\displaystyle \frac{3}{2}}\mu ;1{\displaystyle \frac{\beta ^2}{\alpha ^2}})`$
where $`\alpha =E+m_0`$. The final expression in Eq. (124) completely solves the problem of the generalization of the near threshold expansion technique. For integer values of $`\mu `$ there are no singular Gamma functions (with negative integer argument). Therefore we abolish the regularization and set $`\epsilon =0`$ when using this expression. We have thus found a direct transition from the polarization function (expressed through the integral) to the spectral density in terms of one hypergeometric function for each genuine integral. There is no need to use the recurrence relations available for hypergeometric functions.
As an example we present the (over)simplified case of the two-line diagram with masses $`m`$ and $`m_0m`$ in four space-time dimensions. We cite this example because the expansion of the spectral density and its generalized expansion can be readily compared analytically with the exact result in Eq. (120). Both expansions, the pure and the resummed expansion, will be compared with the exact anaytical result.
The pure expansion of the spectral density near threshold (the second order asymptotic expansion should suffice to show the general features in a short and concise form) is given by
$`\stackrel{~}{\rho }^{das}(E)`$ $`=`$ $`{\displaystyle \frac{\sqrt{2m_0mE}}{8\pi ^2M^{3/2}}}\{1+({\displaystyle \frac{1}{m}}+{\displaystyle \frac{1}{m_0}}{\displaystyle \frac{7}{M}}){\displaystyle \frac{E}{4}}`$ (125)
$`({\displaystyle \frac{1}{m_0^2}}+{\displaystyle \frac{1}{m^2}}+{\displaystyle \frac{12}{m_0m}}{\displaystyle \frac{79}{M^2}}){\displaystyle \frac{E^2}{32}}+O(E^3)\}`$
where $`M=m+m_0`$. As mentioned above, this series breaks down for $`E>2m_0`$ (see Eq. (120)). The analytical expression for the spectral density of the polarization function in Eq. (121) within the generalized asymptotic expansion based on Eq. (124) is given by
$`\stackrel{~}{\rho }^{pas}(E)={\displaystyle \frac{\sqrt{mE(E+2m_0)}}{8\pi ^2(E+M)^{3/2}}}\{{}_{2}{}^{}F_{1}^{}(0,{\displaystyle \frac{1}{2}};{\displaystyle \frac{3}{2}};1{\displaystyle \frac{m_0^2}{(E+m_0)^2}})`$
$`+{\displaystyle \frac{E(E+2m_0)}{8m(E+M)}}{}_{2}{}^{}F_{1}^{}({\displaystyle \frac{1}{2}},1;{\displaystyle \frac{5}{2}};1{\displaystyle \frac{m_0^2}{(E+m_0)^2}})`$
$`{\displaystyle \frac{E^2(E+2m_0)^2}{128m^2(E+M)^2}}(1+{\displaystyle \frac{16m(E+M)}{5(E+m_0)^2}}){}_{2}{}^{}F_{1}^{}(1,{\displaystyle \frac{3}{2}};{\displaystyle \frac{7}{2}};1{\displaystyle \frac{m_0^2}{(E+m_0)^2}})+\mathrm{}\}.`$
We have set the regularization parameter $`\epsilon =0`$ because the spectral density is finite. With $`\epsilon =0`$ the resulting expressions for the hypergeometric functions in Eq. (124) simplify. The first term in the curly brackets of Eq. (4.7) is obviously equal to $`1`$ in this limit because the first parameter of the hypergeometric function vanishes for $`\epsilon =0`$. However, we keep Eq. (4.7) in its given form to show the structure of the contributions. The generalized threshold expansion has the form
$`\stackrel{~}{\rho }^{pas}(E)=g_0(E,m_0)+Eg_1(E,m_0)+E^2g_2(E,m_0)+\mathrm{}`$ (127)
where the functions $`g_j(E,m_0)`$ represent effects of the resummation of small mass effects and are not polynomials in the threshold parameter $`E`$. In the simple one-loop case the hypergeometric functions reduce to elementary functions. For instance,
$`{}_{2}{}^{}F_{1}^{}({\displaystyle \frac{1}{2}},1;{\displaystyle \frac{5}{2}};1{\displaystyle \frac{m_0^2}{(E+m_0)^2}})=`$
$`=`$ $`{\displaystyle \frac{3(E+m_0)}{2E(E+2m_0)}}\left(E+m_0{\displaystyle \frac{m_0^2}{2\sqrt{E(E+2m_0)}}}\mathrm{ln}\left({\displaystyle \frac{E+m_0+\sqrt{E(E+2m_0)}}{E+m_0\sqrt{E(E+2m_0)}}}\right)\right).`$
Higher order contributions are given by hypergeometric functions with larger numerical values of the parameters. They can be simplified by using Gaussian recurrence relations for hypergeometric functions.<sup>2</sup><sup>2</sup>2Note that Eq. (4.7) does not lead to the exact function in Eq. (120) because terms of order $`E^N`$ are missing which originate from the difference part $`\stackrel{~}{\mathrm{\Pi }}^{di}(p)`$ of the correlator. The difference part simply corrects the behaviour of the coefficient functions by the small mass contributions.
In Fig. 11 we compare the exact result of Eq. (112) with the pure expansion in Eq. (125) and the resummed expansion in Eq. (4.7). While the pure expansion breaks down already at the order $`Em_0`$, the convergence of the expansion in Eq. (4.7) breaks down only at $`EM=m+m_0`$. The resummation leads to an essential improvement of the convergence in comparison with the pure threshold expansion. Further examples and their discussion can be found in Ref. .
### 4.8 Large $`s`$ expansion: Expansion of the spectrum at large energies
In the Minkowskian domain where $`s=p^2`$ the situation is closely related to the large $`p`$ expansion if all masses are small compared to the external momentum. We start with a simple example which will be calculated in momentum space as well as in configuration space. In the next paragraph we give more involved examples.
As a first example for an expansion in $`m^2/s`$ we look at a case where the spectral density is given by a finite series only. For the one-loop sunrise-type diagram with one massive and one massless line we can use Feynman parametrization and obtain in (Euclidean) momentum space
$$\stackrel{~}{\mathrm{\Pi }}(p)=\frac{\mathrm{\Gamma }(\epsilon )}{(4\pi )^{2\epsilon }}_0^1(1x)^\epsilon \left(xp^2+m^2\right)^\epsilon ๐x.$$
(129)
In order to determine the spectral desity, the discontinuity can be calculated already on the level of the integrand. A discontinuity appears if $`xp^2+m^2<0`$ which is satisfied for $`p^2<m^2`$. We refer to Appendix C for more details. If we parametrize $`p^2=se^{i\phi }`$, we can reach the cut from both sides $`\phi =\pi `$ and $`\phi =+\pi `$ and can calculate the difference, i.e. the discontinuity
$$\mathrm{Disc}(m^2xs)^\epsilon =\frac{2\pi i}{\mathrm{\Gamma }(\epsilon )\mathrm{\Gamma }(1\epsilon )}(xsm^2)^\epsilon \theta (xsm^2).$$
(130)
For the spectral density we obtain
$$\rho (s)=\frac{1}{2\pi i}\mathrm{Disc}\stackrel{~}{\mathrm{\Pi }}(p)|_{p^2=s}=\frac{1}{(4\pi )^{2\epsilon }\mathrm{\Gamma }(1\epsilon )}_{m^2/s}^1(1x)^\epsilon \left(xsm^2\right)^\epsilon ๐x.$$
(131)
Because the general factor is no longer singular, we can calculate the integral for $`\epsilon =0`$ where the integrand is simply unity. Then we obtain
$$\rho (s)=\frac{1}{(4\pi )^2}_{m^2/s}^1๐x=\frac{1}{(4\pi )^2}\left(1\frac{m^2}{s}\right).$$
(132)
This result can also be obtained by using configuration space techniques. At first sight one may think that this example is too simple to be worth to be calculated using configuration space techniques. Later on we shall see that for the multi-loop case the configuration space method is the only one which leads to an analytical result.
In configuration space the spectral density of the one-loop diagram with one massive and one massless line in four-dimensional spacetime is given by Eq. (112). The integral contains the modified Bessel function of the first kind, $`I_1(z)`$. Because $`K_1(x)`$ is a pure exponentially decreasing function in its asymptotic expansion while $`K_1(x)\pm i\pi I(x)`$ increases exponentially along the real axis if continued in the upper resp. lower complex plane (note the ambiguity due to Stokesโ phenomenon mentioned in Appendix C), we can divide according to
$`i\pi {\displaystyle _ฯต^{ฯต+i\mathrm{}}}I_1(z)f(z)๐z`$ $`=`$ $`{\displaystyle _ฯต^{ฯต+i\mathrm{}}}(i\pi I_1(z)K_1(z))f(z)๐z+{\displaystyle _ฯต^{ฯต+i\mathrm{}}}K_1(z)๐z,`$ (133)
$`i\pi {\displaystyle _{ฯตi\mathrm{}}^ฯต}I_1(z)f(z)๐z`$ $`=`$ $`{\displaystyle _{ฯตi\mathrm{}}^ฯต}(i\pi I_1(z)+K_1(z))f(z)๐z{\displaystyle _{ฯตi\mathrm{}}^ฯต}K_1(z)๐z`$ (134)
and can lower the paths to the negative resp. positive real axis by using Cauchyโs theorem realizing that the quarter circle integrals in the corresponding quadrant vanish (see Ref. ). The origin and the negative real axis is circumvented by two half-circle paths $`C_{}`$ (lower half plane) and $`C_+`$ (upper half plane) in positive direction. The result found in Ref. reads
$`i\pi {\displaystyle _{ci\mathrm{}}^{c+i\mathrm{}}}I_1(x\sqrt{s})\mathrm{\Pi }(x)x^2๐x={\displaystyle _ฯต^{\mathrm{}}}K_1(r\sqrt{s})\left(2\mathrm{\Pi }(r)\mathrm{\Pi }(e^{i\pi }r)\mathrm{\Pi }(e^{i\pi }r)\right)r^2๐r`$
$`+{\displaystyle _C_{}}(i\pi I_1(z\sqrt{s})+K_1(z\sqrt{s}))\mathrm{\Pi }(z)z^2๐z+{\displaystyle _{C_+}}(i\pi I_1(z\sqrt{s})K_1(z\sqrt{s}))\mathrm{\Pi }(z)z^2๐z.`$
For our simple example we have to take the correlator function $`\mathrm{\Pi }(x)=D(x,m)D(x,0)`$. While the massless propagator is invariant under the multiplication of $`e^{\pm i\pi }=\pm 1`$ of its argument, i.e. $`D(x,0)=D(x,0)`$, for the massive propagator we obtain
$$D(e^{\pm i\pi }r,m)=\frac{(mr)}{(2\pi )^2r^2}\left(K_1(mr)\pm i\pi I_1(mr)\right)$$
(136)
where $`K_\lambda (e^{\pm i\pi }x)=e^{i\pi \lambda }K_\lambda (x)i\pi I_\lambda (x)`$ is used. We obtain $`2\mathrm{\Pi }(r)\mathrm{\Pi }(e^{i\pi }r)\mathrm{\Pi }(e^{i\pi }r)=0`$ which results in the vanishing of the first part in Eq. (4.8). We are therefore left with the two semicircle integrals which we can combine in a different manner
$$\rho (s)=\frac{2\pi i}{\sqrt{s}}_{C_{}+C_+}I_1(z\sqrt{s})\mathrm{\Pi }(z)z^2๐z\frac{2}{\sqrt{s}}_{C_{}C_+}K_1(z\sqrt{s})\mathrm{\Pi }(z)z^2๐z.$$
(137)
In expanding the first integrand for small values of $`r=|z|`$ we obtain
$$I_1(z\sqrt{s})\mathrm{\Pi }(z)z^2=\frac{\sqrt{s}}{2(2\pi )^4z}+O(z),$$
(138)
the integral results in
$$_{C_{}+C_+}I_1(z\sqrt{s})\mathrm{\Pi }(z)z^2๐z=\frac{\sqrt{s}}{2(2\pi )^4}\left[\mathrm{ln}r+i\phi \right]_\pi ^\pi =\frac{i\sqrt{s}}{16\pi ^3}.$$
(139)
For the second integrand the expansion results in many more terms,
$`K_1(z\sqrt{s})\mathrm{\Pi }(z)z^2`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^4z^3\sqrt{s}}}+{\displaystyle \frac{1}{4(2\pi )^4z\sqrt{s}}}(m^2\mathrm{ln}(m^2)+s\mathrm{ln}s)`$ (140)
$`+{\displaystyle \frac{m^2+s}{2(2\pi )^4z\sqrt{s}}}\left(\mathrm{ln}\left({\displaystyle \frac{z}{2}}\right)+\gamma _E{\displaystyle \frac{1}{2}}\right).`$
However one has
$`{\displaystyle _{C_{}C_+}}{\displaystyle \frac{dz}{z^3}}`$ $`=`$ $`\left[{\displaystyle \frac{1}{2z^2}}\right]_{C_{}C_+}=\left[{\displaystyle \frac{1}{2r^2e^{2i\phi }}}\right]_{\phi =\pi }^0\left[{\displaystyle \frac{1}{2r^2e^{2i\phi }}}\right]_{\phi =0}^\pi =0,`$
$`{\displaystyle _{C_{}C_+}}{\displaystyle \frac{dz}{z}}`$ $`=`$ $`\left[\mathrm{ln}z\right]_{C_{}C_+}=\left[\mathrm{ln}ฯต+i\phi \right]_{\phi =\pi }^0\left[\mathrm{ln}ฯต+i\phi \right]_{\phi =0}^\pi =\pi \pi =0.`$ (141)
The only non-vanishing contribution is given by
$$_{C_{}C_+}\frac{dz}{z}\mathrm{ln}z=\frac{1}{2}\left[\mathrm{ln}^2z\right]_{C_{}C_+}=\pi ^2.$$
(142)
Therefore, we obtain
$$_{C_{}C_+}K_1(z\sqrt{s})\mathrm{\Pi }(z)z^2๐z=\frac{m^2+s}{32\pi ^2\sqrt{s}}.$$
(143)
The spectral density, finally, reads again
$$\rho (s)=\frac{2\pi i}{\sqrt{s}}\frac{i\sqrt{s}}{16\pi ^3}\frac{2}{\sqrt{s}}\frac{m^2+s}{32\pi ^2\sqrt{s}}=\frac{1}{16\pi ^2}\left(1\frac{m^2}{s}\right)$$
(144)
as in Eq. (132).
### 4.9 Large $`s`$ expansion for massive one- and multiloop diagrams
A second example that we treat is the one-loop diagram with two different masses $`m_1`$ and $`m_2`$. We start again with the calculation in momentum space to obtain
$$\stackrel{~}{\mathrm{\Pi }}(p)=\frac{\mathrm{\Gamma }(2D/2)}{(4\pi )^{D/2}}_0^1\left(x(1x)p^2+(1x)m_1^2+xm_2^2\right)^{D/22}๐x.$$
(145)
Again the discontinuity can be calculated already on the level of integrands. Using Eq. (C15) we obtain
$`\mathrm{Disc}\left(x(1x)s+(1x)m_1^2+xm_2^2\right)^\epsilon =`$
$`=`$ $`{\displaystyle \frac{2\pi i}{\mathrm{\Gamma }(\epsilon )\mathrm{\Gamma }(1\epsilon )}}\left(x(1x)s(1x)m_1^2xm_2^2\right)^\epsilon \theta \left(x(1x)s(1x)m_1^2xm_2^2\right).`$
In order to write down the spectral density we have to determine the limits for $`x`$ given by the positiveness of the theta function. The zeros of the argument of the theta function can be calculated to be
$$x_{1,2}=\frac{1}{2}\left(1+\frac{m_1^2}{s}\frac{m_2^2}{s}\right)\pm \frac{1}{2}\sqrt{\lambda (1,m_1^2/s,m_2^2/s)}$$
(147)
where
$$\lambda (x,y,z)=x^2+y^2+z^22xy2xz2yz$$
(148)
is Kรคllรฉnโs lambda function. In terms of these two zeros $`x_1x_2`$ we can write the condition for the argument of the theta function as
$$(xx_1)(x_2x)s0x_1xx_2.$$
(149)
The two zeros take real values if
$$s^2\lambda (1,m_1^2/s,m_2^2/s)=\left(s(m_1m_2)^2\right)\left(s(m_1+m_2)^2\right)0.$$
(150)
This is the case for $`s(m_1+m_2)^2`$ or $`s(m_1m_2)^2`$. We do not consider the part of the spectral density up to the pseudothreshold $`(m_1m_2)^2`$ because in this case both $`x_1`$ and $`x_2`$ are negative. Starting from the threshold $`(m_1+m_2)^2`$, both $`x_1`$ and $`x_2`$ take values between $`0`$ and $`1`$. We obtain
$`\rho (s)`$ $`=`$ $`{\displaystyle \frac{1}{2\pi i}}\mathrm{Disc}\stackrel{~}{\mathrm{\Pi }}(s)={\displaystyle \frac{1}{(4\pi )^{2\epsilon }\mathrm{\Gamma }(1\epsilon )}}{\displaystyle _{x_1}^{x_2}}s^\epsilon (xx_1)^\epsilon (x_2x)^\epsilon ๐x`$ (151)
$`=`$ $`{\displaystyle \frac{s^\epsilon \mathrm{\Gamma }(1\epsilon )}{(4\pi )^{2\epsilon }\mathrm{\Gamma }(22\epsilon )}}\sqrt{\lambda (1,m_1^2/s,m_2^2/s)}`$
where we have used the substitution $`x^{}=(xx_1)/(x_2x_1)`$ and Eulerโs Beta function
$$B(\alpha ,\beta )=\frac{\mathrm{\Gamma }(\alpha )\mathrm{\Gamma }(\beta )}{\mathrm{\Gamma }(\alpha +\beta )}=_0^1x^{\alpha 1}(1x)^{\beta 1}๐x.$$
(152)
For $`D=4`$, i.e. $`\epsilon =0`$, we obtain
$$\rho (s)=\frac{1}{(4\pi )^2}\sqrt{12\left(\frac{m_1^2}{s}+\frac{m_2^2}{s}\right)+\left(\frac{m_1^2}{s}\frac{m_2^2}{s}\right)^2}.$$
(153)
For the case $`m_2=0`$ we recover the spectral density of the previous example. If we consider an expansion in terms of large values of $`s`$ compared to the threshold value $`s=(m_1+m_2)^2`$, i.e. if we expand in small values of $`z`$ for $`z=(m_1+m_2)^2/s`$, we obtain
$$\rho (s)=\frac{1}{(4\pi )^2}\left(1\frac{m_1^2+m_2^2}{s}+O\left(\frac{1}{s^2}\right)\right).$$
(154)
The same result shall now be calculated in a different way. Consider the correlator function
$$\stackrel{~}{\mathrm{\Pi }}(p)=2\pi ^{\lambda +1}_0^{\mathrm{}}\left(\frac{px}{2}\right)^\lambda J_\lambda (px)D(x,m_1)D(x,m_2)x^{2\lambda +1}๐x.$$
(155)
We can expand the propagators $`D(x,m_i)`$ in $`m_i/\sqrt{s}`$, leading to pure powers in $`x`$. Starting from Eq. (86), the integration of the Bessel function of the first kind with (non-integer) powers of $`x`$ can be cast into a more appropriate form
$$_0^{\mathrm{}}x^\mu \left(\frac{px}{2}\right)^\lambda J_\lambda (px)x^{2\lambda +1}๐x=\left(\frac{p}{2}\right)^{2\lambda 2\mu }\frac{\mathrm{\Gamma }(\lambda +1+\mu /2)}{2\mathrm{\Gamma }(\mu /2)}.$$
(156)
We can then immediately calculate the discontinuity and obtain
$$\frac{1}{2\pi i}\mathrm{Disc}_0^{\mathrm{}}x^\mu \left(\frac{px}{2}\right)^\lambda J_\lambda (px)x^{2\lambda +1}๐x=\frac{(s/4)^{\lambda 1\mu /2}}{2\mathrm{\Gamma }(\lambda \mu /2)\mathrm{\Gamma }(\mu /2)}.$$
(157)
Using the expansion of the Bessel functions $`K_\lambda (z)`$ at small values of $`z`$, the expansion of the propagator for $`\lambda `$ close to $`\lambda =1`$ for small mass $`m`$ (compared to $`\sqrt{s}`$) reads
$`D(x,m)={\displaystyle \frac{(mx)^\lambda K_\lambda (mx)}{(2\pi )^{\lambda +1}x^{2\lambda }}}`$ (158)
$`=`$ $`{\displaystyle \frac{2^{\lambda 1}\mathrm{\Gamma }(\lambda )}{(2\pi )^{\lambda +1}x^{2\lambda }}}\left(1{\displaystyle \frac{1}{1\lambda }}\left({\displaystyle \frac{mx}{2}}\right)^2{\displaystyle \frac{\mathrm{\Gamma }(1\lambda )}{\mathrm{\Gamma }(1+\lambda )}}\left({\displaystyle \frac{mx}{2}}\right)^{2\lambda }+O((mx)^4,(mx)^{2+2\lambda })\right).`$
Eq. (155) contains two of these propagators. We now expand up to the power $`m^2`$ and keep $`\lambda `$ close to $`1`$. We obtain a zeroth order contribution, two contributions of the order $`m^2`$, and two contributions of the order $`m^{2(1+\lambda )}`$. We consider these contributions in turn.
* zeroth order ($`\mu =4\lambda `$):
$$\frac{2\pi ^{\lambda +1}(s/4)^{\lambda 1}}{2\mathrm{\Gamma }(2\lambda )\mathrm{\Gamma }(\lambda )}\left(\frac{2^{\lambda 1}\mathrm{\Gamma }(\lambda )}{(2\pi )^{\lambda +1}}\right)^2\frac{2\pi ^2}{2\mathrm{\Gamma }(2)\mathrm{\Gamma }(1)}\left(\frac{\mathrm{\Gamma }(1)}{(2\pi )^2}\right)^2=\frac{1}{16\pi ^2}$$
(159)
* second order ($`\mu =24\lambda `$, $`i=1,2`$):
$`{\displaystyle \frac{2\pi ^{\lambda +1}(s/4)^{\lambda 2}}{2\mathrm{\Gamma }(2\lambda 1)\mathrm{\Gamma }(\lambda 1)}}\left({\displaystyle \frac{2^{\lambda 1}\mathrm{\Gamma }(\lambda )}{(2\pi )^{\lambda +1}}}\right)^2{\displaystyle \frac{(m_i^2/4)}{1\lambda }}\left(\mathrm{\Gamma }(\lambda )=(\lambda 1)\mathrm{\Gamma }(\lambda 1)\right)`$
$`=`$ $`{\displaystyle \frac{2\pi ^{\lambda +1}(s/4)^{\lambda 2}}{2\mathrm{\Gamma }(2\lambda 1)\mathrm{\Gamma }(\lambda )}}\left({\displaystyle \frac{2^{\lambda 1}\mathrm{\Gamma }(\lambda )}{(2\pi )^{\lambda +1}}}\right)^2{\displaystyle \frac{m_i^2}{4}}{\displaystyle \frac{2\pi ^2(s/4)^1}{2\mathrm{\Gamma }(1)\mathrm{\Gamma }(1)}}\left({\displaystyle \frac{\mathrm{\Gamma }(1)}{(2\pi )^{\lambda +1}}}\right)^2{\displaystyle \frac{m_i^2}{4}}={\displaystyle \frac{m_i^2}{16\pi ^2s}}`$
* order $`m^{2(1+\lambda )}`$ (also called โsecond order primedโ, $`\mu =2\lambda `$, $`i=1,2`$):
$$\frac{(s/4)^1}{\mathrm{\Gamma }(\lambda \mu /2)\mathrm{\Gamma }(\mu /2)}\left(\frac{2^{\lambda 1}\mathrm{\Gamma }(\lambda )}{(2\pi )^{\lambda +1}}\right)^2\frac{2\mathrm{\Gamma }(1\lambda )}{\mathrm{\Gamma }(1+\lambda )}\left(\frac{m_i^2}{4}\right)^\lambda .$$
(161)
Since the third contribution vanishes for $`\mu 2\lambda `$, only the first two contributions have to be taken into account. One ends up with
$$\rho (s)=\frac{1}{16\pi ^2}\left(1\frac{m_1^2}{s}\frac{m_2^2}{s}+O\left(\frac{1}{s^2}\right)\right)$$
(162)
in agreement with the result obtained earlier. Moreover, we can predict the result for a $`n`$-loop sunrise-type diagram with different masses $`m_1,\mathrm{}m_{n+1}`$ in four space-time dimensions where $`\mu =2(n+1)\lambda `$ (zeroth order contribution), $`\mu =22(n+1)\lambda `$ (second order contribution), and $`\mu =2n\lambda `$ (second order primed contribution which does not vanish for $`n>1`$),
$`\rho (s)`$ $`=`$ $`2\pi ^{\lambda +1}\left({\displaystyle \frac{2^{\lambda 1}\mathrm{\Gamma }(\lambda )}{(2\pi )^{\lambda +1}}}\right)^{n+1}{\displaystyle \frac{(s/4)^{n\lambda 1}}{2\mathrm{\Gamma }((n+1)\lambda )\mathrm{\Gamma }(n\lambda )}}(1+`$
$`+{\displaystyle \frac{((n+1)\lambda 1)(n\lambda 1)}{1\lambda }}{\displaystyle \underset{i=1}{\overset{n+1}{}}}{\displaystyle \frac{m_i^2}{s}}{\displaystyle \frac{\mathrm{\Gamma }((n+1)\lambda )\mathrm{\Gamma }(1\lambda )}{\mathrm{\Gamma }((n1)\lambda )\mathrm{\Gamma }(1+\lambda )}}{\displaystyle \underset{i=1}{\overset{n+1}{}}}\left({\displaystyle \frac{m_i^2}{s}}\right)^\lambda +O\left({\displaystyle \frac{1}{s^2}}\right)).`$
The second and third term have to be expanded in $`\epsilon =1\lambda `$ in order that the singularity in $`\epsilon `$ cancels. After a few simplifications using expecially
$$\psi (n+1)=\psi (n)+\frac{1}{n}$$
(164)
for the polygamma function, we finally obtain for $`\lambda =1`$
$$\rho (s)=\frac{s^{n1}}{(4\pi )^{2n}n!(n1)!}\left(1+n\underset{i=1}{\overset{n+1}{}}\left((n1)\left(\mathrm{ln}\left(\frac{m_i^2}{s}\right)+2(\psi (n)+\gamma _E)\right)n\right)\frac{m_i^2}{s}+O\left(\frac{1}{s^2}\right)\right).$$
(165)
This result is also valid for $`n=1`$ and again confirms our previous results. Note that $`(\psi (n)+\gamma _E)`$ is a rational number because Eulerโs constant cancels out. Indeed, we have
$$\psi (n)+\gamma _E=\underset{k=1}{\overset{n1}{}}\frac{1}{k}.$$
(166)
Note, however, that for $`n>1`$ logarithmic contributions $`\mathrm{ln}(m_i^2/s)`$ appear in the spectral density.
## 5 Non-standard propagators and other exotic settings
In this section we deal with modifications of the standard propagators in sunrise-type diagrams in any space-time dimension. First we consider odd-dimensional space-time which in most of the cases allows one to calculate the sunrise-type integrals in closed form. Returning to space-time dimensions close to four, we then deal with larger powers of propagators leading to larger values for the Bessel function indices. In this case one can apply recurrence relations to reduce integrals containing higher powers of the propagators to a set of master integrals . The master integrals themselves show Laplace-type asymptotics (for a definition see e.g. Ref. ). Finally we deal with nontrivial numerators in two examples involving vacuum bubbles.
### 5.1 Odd-dimensional case
Compared to four-dimensional space-time, integrals containing products of Bessel functions in odd dimensional space-time are analytically solvable because the Bessel functions simplify significantly. It is interesting to note that the evaluation of Eq. (18) can be done in a closed form for any number of internal lines in odd-dimensional space-time. One might suspect that this special case is too simple to have any practical applications. However, at the end of this subsection we will refer to many applications from various fields of physics involving odd-dimensional space-time. As the simplest example we take three-dimensional space-time $`DD_0=3`$. For $`D_0=3`$ (and $`\lambda _0=(D_02)/2=1/2`$) the propagator in Eq. (7) reads
$$D(x,m)D_3(x,m)=\frac{\sqrt{mx}K_{1/2}(mx)}{(2\pi )^{3/2}x}=\frac{e^{mx}}{4\pi x}.$$
(167)
With $`\lambda =\lambda _0=1/2`$ the weight function becomes
$$\left(\frac{px}{2}\right)^{1/2}J_{1/2}(px)=\frac{2}{\sqrt{\pi }}\frac{\mathrm{sin}(px)}{px}$$
(168)
after angular integration. The explicit result for the $`(n+1)`$-line sunrise diagram is then given by the integral
$`\stackrel{~}{\mathrm{\Pi }}(p)`$ $`=`$ $`4\pi {\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\mathrm{sin}(px)}{px}}{\displaystyle \frac{e^{Mx}}{(4\pi x)^{n1}}}(\mu x)^{2ฯต}๐x`$ (169)
$`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(1n+2ฯต)}{2ip(4\pi )^n}}\left[(Mip)^{n12ฯต}(M+ip)^{n12ฯต}\right]\mu ^{2ฯต}`$
where $`ฯต`$ is used for regularization and $`M=m_i`$. Note that we have used an unorthodox regularization method by multiplying a factor $`(\mu x)^{2\epsilon }`$. We will return to this point later on.
We consider some particular cases of Eq. (169) for different values of $`n`$. For $`n=0`$ we simply recover the (Euclidean) propagator function $`\stackrel{~}{\mathrm{\Pi }}(p)=(M^2+p^2)^1`$ with the discontinuity
$$\rho (s)=\frac{\text{Disc }\stackrel{~}{\mathrm{\Pi }}(p)}{2\pi i}=\frac{1}{2\pi i}\left(\stackrel{~}{\mathrm{\Pi }}(p)|_{p^2=s\mathrm{exp}(i\pi )}\stackrel{~}{\mathrm{\Pi }}(p)|_{p^2=s\mathrm{exp}(i\pi )}\right)=\delta (sm^2)$$
(170)
where $`s`$ is the squared energy, $`s=p^2`$. As remarked on earlier it is appropriate to call this expression the spectral density associated with the diagram. For $`n=1`$ the answer for the polarization function $`\stackrel{~}{\mathrm{\Pi }}(p)`$ is still finite (no regularization is required) and is given by
$$\stackrel{~}{\mathrm{\Pi }}(p)=\frac{1}{8\pi ip}\mathrm{ln}\left(\frac{M+ip}{Mip}\right).$$
(171)
The spectral density, i.e. the discontinuity of Eq. (171) across the cut in the complex $`p^2`$-plane is given by
$$\rho (s)=\frac{1}{8\pi \sqrt{s}}\theta (sM^2),s=p^2,s>0$$
(172)
which is nothing but the three-dimensional two-particle phase space. This can be immediately checked by direct computation. The cases with $`n>1`$ have more structure and therefore are more interesting. For the genuine sunrise diagram with $`n=2`$, Eq. (169) leads to
$$\stackrel{~}{\mathrm{\Pi }}(p)=\frac{1}{32\pi ^2}\left(\frac{1}{ฯต}\frac{M}{ip}\mathrm{ln}\left(\frac{M+ip}{Mip}\right)\mathrm{ln}\left(\frac{M^2+p^2}{\mu ^2}\right)\right).$$
(173)
While in Eq. (173) the arbitrary scale $`\mu ^2`$ appears due to regularization, the spectral density
$$\rho (s)=\frac{\sqrt{s}M}{32\pi ^2\sqrt{s}}\theta (sM^2)$$
(174)
is independent of this scale. This is because the spectral density is again finite and, therefore, independent of the regularization used. The general formula for the spectral density for any $`n>0`$ in $`D=3`$ can be extracted from Eq. (169). It reads
$$\rho (s)=\frac{(\sqrt{s}M)^{n1}}{2(4\pi )^n(n1)!\sqrt{s}}\theta (sM^2).$$
(175)
We now want to comment on the relation between the momentum subtraction and our unorthodox dimensional regularization. Taking Eq. (169) for $`n=2`$ with momentum subtraction at the origin, one obtains
$$\stackrel{~}{\mathrm{\Pi }}(p)=_0^{\mathrm{}}\left(\frac{\mathrm{sin}(px)}{px}1\right)\frac{e^{Mx}}{(4\pi )^2x}(\mu x)^{2ฯต}๐x$$
(176)
which is UV-finite even for $`ฯต=0`$ because there is no singularity at the origin. For practical computations it is convenient to keep the factor $`(\mu x)^{2ฯต}`$ in the integrand since this factor gives a meaning to each of the two terms in the round brackets in Eq. (176). Using this factor we obtain
$`\stackrel{~}{\mathrm{\Pi }}(p)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}\left({\displaystyle \frac{\mathrm{sin}(px)}{px}}1\right){\displaystyle \frac{e^{Mx}}{(4\pi )^2x}}(\mu ^2x^2)^ฯต๐x`$ (177)
$`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(1+2ฯต)}{2ip(4\pi )^2}}\left[(Mip)^{12ฯต}(M+ip)^{12ฯต}\right]\mu ^{2ฯต}{\displaystyle \frac{\mathrm{\Gamma }(2ฯต)}{(4\pi )^2}}\left({\displaystyle \frac{\mu }{M}}\right)^{2ฯต}`$
$`=`$ $`{\displaystyle \frac{1}{32\pi ^2}}\left\{{\displaystyle \frac{M}{ip}}\mathrm{ln}\left({\displaystyle \frac{M+ip}{Mip}}\right)+\mathrm{ln}\left({\displaystyle \frac{M^2+p^2}{M^2}}\right)\right\}.`$
The poles cancel in this expression and the arbitrary scale $`\mu `$ changes to $`M`$. This corresponds to a transition from MS-type renormalization schemes to the momentum subtraction scheme (in this particular case with a subtraction at the origin). Since the spectral density $`\rho (s)`$ is finite it can be computed using any regularization scheme as can be seen from Eqs. (173) and (177).
We mention that in the three-dimensional case the spectral density $`\rho (s)`$ can also be found for general values of $`n`$ by traditional methods since the three-dimensional case is sufficiently simple. By using the convolution for the evaluation of spectral densities one stays in the same class of functions, i.e. polynomials in the variable $`\sqrt{s}`$ divided by $`\sqrt{s}`$. The general form of the convolution equation in $`D`$-dimensional space-time reads
$$\mathrm{\Phi }_n(s)=\mathrm{\Phi }_k(s_1)\mathrm{\Phi }_p(s_2)\mathrm{\Phi }_1(s,s_1,s_2)๐s_1๐s_2,k+p+1=n.$$
(178)
For the particular case of three-dimensional space-time the kernel $`\mathrm{\Phi }_1(p^2,m_1^2,m_2^2)`$ is given by
$$(2\pi )^2\mathrm{\Phi }_1(p^2,m_1^2,m_2^2)=\delta (k^2m_1^2)\delta ((pk)^2m_2^2)d^3k$$
(179)
or explicitly by
$$\mathrm{\Phi }_1(s,s_1,s_2)=\frac{1}{8\pi \sqrt{s}}\theta (s(\sqrt{s_1}+\sqrt{s_2})^2).$$
(180)
Eq. (180) can be seen to be the two-particle phase space in three dimensions (cf. Eq. (172)). This is a rather simple example. However, our technique retains its efficiency for large $`n`$.
We list some potential applications of the general results obtained in this paragraph for odd-dimensional space-time. In three space-time dimensions our results can be used to compute phase space integrals for particles in jets where the momentum along the direction of the jet is fixed . Another application can be found in phase transitions, for instance the three-dimensional QCD which emerges as the high temperature limit of the ordinary theory of strong interactions for the quark-gluon plasma (see e.g. ). Three-dimensional models are also used to study the question of dynamical mass generation and the infrared structure of the models of Quantum Field Theory in general and some problems of QCD .
Note that particular models of different space-time dimensions are very useful because their properties may be simpler and may thus allow one to study general features of the underlying field theory. For example, in six-dimensional space-time the simplest model of quantum field theory $`\varphi ^3`$ is asymptotically free and can be used for simulations of some features of QCD .
As we have already mentioned before, closed form analytical results can be obtained for odd dimensional space-time $`D=3,5,7,\mathrm{}`$. For example, the propagator in five-dimensional space-time ($`\lambda =3/2`$) reads
$$D(x,m)D_5(x,m)=\frac{(mx)^{3/2}K_{3/2}(mx)}{(2\pi )^{5/2}x^3}=\frac{e^{mx}}{8\pi ^2x^3}(1+mx)$$
(181)
which assures that the basic integration can be performed in terms of elementary functions (powers and logarithms) again.
Applications for odd space-time dimensions other than three can be found in some models in unified field theories, or also in a general analysis of the divergence structure in QFT. Five-dimensional models of QFT are rather popular for general purposes . They have useful applications for Yang-Mills theories in five-dimensional space-time where the UV structure of the models can be analyzed .
Finally, non-standard space-time dimensions are useful in order to obtain estimates for the standard case $`D_0=4`$. The feature that $`K_\nu (x)>K_\mu (x)`$ holds for modified Bessel functions of the second kind when $`\nu >\mu `$ and for positive $`x`$-arguments leads to $`K_{3/2}(x)>K_1(x)>K_{1/2}(x)`$ which can be used for stringent numerical estimates for integrals containing these Bessel functions.
We conclude this paragraph by noting that the $`x`$-space techniques allow one to compute sunrise-type diagrams in closed form in terms of elementary functions as long as one is dealing with odd-dimensional space-times. The resulting expressions are rather simple and can be directly used for applications. Having the complete formulas at hand, there is no need to expand in the parameters of the diagram such as masses or external momenta.
### 5.2 Large powers of propagators
In many physics applications one encounters large powers of propagators. Recent examples are the calculation of the corrections to $`B^0\overline{B}^0`$ mixing in perturbative QCD (see Fig. 12) or the large mass expansion for the contribution of charged scalars to the muon anomalous magnetic moment .
When calculating moments of the spectral density of such a diagram by using packages like MATAD for the automatic calculation of Feynman diagrams , high powers of propagators are generated. Another example are sunrise-type diagrams with external lines at small momenta $`q_i`$ (see Fig. 13). These diagrams appear in calculations when there are weak external fields as, for example, in sum rule applications or in special cases of high precision calculations in ChPT .
In momentum space larger powers of propagators are treated through recurrence relations based on the integration-by-parts technique . In configuration space larger powers of the propagators give rise to larger values of the indices for the corresponding Bessel functions. We now describe a technique using recurrence relations which allows one to reduce large indices of Bessel functions .
Note that the direct reduction of a sunrise-type diagram to a standard set of master integrals with the help of algebraic computer systems is rather time-consuming with present momentum space techniques. In practice the computation proceeds through the use of a table of integrals with given powers of the denominators. One would have to set up a three-dimensional table for a given total power $`N`$. The number of entries (even when accounting for the appropriate symmetries) still grows as fast as $`N^3`$ which is large for the large values of $`N`$ needed in some present applications. Within our method one first re-expresses the relevant integrals through a one-parameter set of integrals which are then solved explicitly. For large $`N`$ the number of entries increases as a first power of $`N`$ (the number of elements for the $`I_0(q)`$ basis is given by $`2[N/2]5`$ where $`[z]`$ is an integer part of $`z`$) which considerably reduces the time consumption in a computer evaluation.
In order to show the applicability of our method we shall calculate three examples that were solved before with the help of momentum space techniques . In Ref. Broadhurst considered general three-loop bubbles $`B_N`$. A subclass of these are the sunrise-type three-loop bubbles. In momentum space they read
$`B_N(0,0,n_3,n_4,n_5,n_6)={\displaystyle }{\displaystyle \frac{d^Dkd^Dld^Dp}{m^{3D}(\pi ^{D/2}\mathrm{\Gamma }(3D/2))^3}}\times `$ (182)
$`\times {\displaystyle \frac{m^{2n_3}}{((p+k)^2+m^2)^{n_3}}}{\displaystyle \frac{m^{2n_4}}{((p+l)^2+m^2)^{n_4}}}{\displaystyle \frac{m^{2n_5}}{((p+k+l)^2+m^2)^{n_5}}}{\displaystyle \frac{m^{2n_6}}{(p^2+m^2)^{n_6}}}`$
with two propagators absent ($`n_1=n_2=0`$). Actually we choose the remaining indices $`n_3`$, $`n_4`$, $`n_5`$, and $`n_6`$ such that the results become finite. We shall not only calculate the finite part but also the part porportional to $`\epsilon `$ in order to be able to compare with . Written in configuration space, the particular subset of bubble diagrams $`B_N`$ is given by
$`B_N(0,0,n_3,n_4,n_5,n_6)={\displaystyle \frac{2(64\pi ^4)^{2\epsilon }}{(\mathrm{\Gamma }(1+\epsilon ))^3\mathrm{\Gamma }(2\epsilon )}}m^{2(n_3+n_4+n_5+n_6)12+6\epsilon }\times `$ (183)
$`\times {\displaystyle _0^{\mathrm{}}}D^{(n_31)}(x,m)D^{(n_41)}(x,m)D^{(n_51)}(x,m)D^{(n_61)}(x,m)x^{2\lambda +1}dx.`$
It is clear that one will end up with integrals of products of four Bessel functions with non-integer indices and a non-integer power of $`x`$. We first discuss three examples in the next three subsections which will be followed by more general considerations on a general reduction procedure which allows one to reduce all integrals containing products of Bessel functions to the following set of two master integrals,
$$L_4(r):=_0^{\mathrm{}}\left(K_0(\xi )\right)^4\xi ^r๐\xi \text{and}L_4^l(r):=_0^{\mathrm{}}\left(K_0(x)\right)^4\xi ^r\mathrm{ln}(e^{\gamma _E}\xi /2)๐x$$
(184)
where the index โ$`l`$โ in $`L_4^l(r)`$ reflects the logarithm appearing in the integrand. We will then add some considerations on the basic integrals $`L_4(r)`$ and $`L_4^l(r)`$.
### 5.3 The example $`B_N(0,0,2,2,2,2)`$
We start with the example which is represented by the diagram in Fig. 14(b). Each of the lines is modified (indicated by the dots on the lines) which means that instead of the propagators $`D(x,m)`$ we have to use
$$D^{(1)}(x,m)=\frac{d^Dp}{(2\pi )^D}\frac{e^{ip_\mu x^\mu }}{(p^2+m^2)^2}=\frac{(x/m)^{1\lambda }}{2(2\pi )^{\lambda +1}}K_{\lambda 1}(mx)=\frac{(x/m)^\epsilon }{2(2\pi )^{2\epsilon }}K_\epsilon (mx).$$
(185)
One obtains ($`\xi =mx`$)
$$B_N(0,0,2,2,2,2)=\frac{2^{12\epsilon }}{(\mathrm{\Gamma }(1+\epsilon ))^3\mathrm{\Gamma }(2\epsilon )}_0^{\mathrm{}}\left(K_\epsilon (\xi )\right)^4\xi ^{3+2\epsilon }๐\xi .$$
(186)
We now can use the general formula
$$\left[\frac{K_\nu (z)}{\nu }\right]_{\nu =\pm n}=\pm \frac{1}{2}n!\underset{k=0}{\overset{n1}{}}\left(\frac{z}{2}\right)^{kn}\frac{K_k(z)}{k!(nk)},n\{0,1,\mathrm{}\}$$
(187)
to expand the Bessel function in a series with respect to its index. In case of $`K_\epsilon (z)`$, however, the first derivative vanishes and we obtain $`K_\epsilon (z)=K_0(z)+O(\epsilon ^2)`$. Therefore, in expanding
$`{\displaystyle \frac{2^{12\epsilon }\left(K_\epsilon (\xi )\right)^4\xi ^{3+2\epsilon }}{(\mathrm{\Gamma }(1+\epsilon ))^3\mathrm{\Gamma }(2\epsilon )}}`$ $`=`$ $`2\left(K_0(\xi )\right)^4\xi ^3\left(1+(1+2\gamma _E2\mathrm{ln}2+2\mathrm{ln}\xi )\epsilon +O(\epsilon ^2)\right)=`$ (188)
$`=`$ $`2\left(K_0(\xi )\right)^4\xi ^3\left(1+\left(1+2\mathrm{ln}(e^{\gamma _E}\xi /2)\right)\epsilon +O(\epsilon ^2)\right)`$
we obtain
$`B_N(0,0,2,2,2,2)`$ $`=`$ $`2(1+\epsilon ){\displaystyle _0^{\mathrm{}}}\left(K_0(\xi )\right)^4\xi ^3๐\xi +4\epsilon {\displaystyle _0^{\mathrm{}}}\left(K_0(\xi )\right)^4\xi ^3\mathrm{ln}(e^{\gamma _E}\xi /2)๐\xi =`$ (189)
$`=`$ $`2(1+\epsilon )L_4(3)+4\epsilon L_4^l(3).`$
We have compared our numerical result with the analytical result in and have found agreement. In fact, from the transcendality structure of the results in one can numerically obtain the appropriate coefficients multiplying the transcendentals. One finds
$`L_4(3)`$ $`=`$ $`{\displaystyle \frac{3}{16}}+{\displaystyle \frac{7}{32}}\zeta (3),`$
$`L_4^l(3)`$ $`=`$ $`{\displaystyle \frac{3}{32}}+{\displaystyle \frac{3}{4}}\mathrm{Li}_4\left({\displaystyle \frac{1}{2}}\right){\displaystyle \frac{17\pi ^4}{1920}}{\displaystyle \frac{\pi ^2}{32}}\left(\mathrm{ln}2\right)^2+{\displaystyle \frac{1}{32}}\left(\mathrm{ln}2\right)^4+{\displaystyle \frac{49}{128}}\zeta (3).`$ (190)
### 5.4 The example $`B_N(0,0,2,2,2,1)`$
In the diagram of Fig. 14(a) one of the lines is not modified. Therefore, we have to deal with one regular propagator factor
$$D^{(0)}(x,m)=D(x,m)=\frac{(x/m)^\lambda }{(2\pi )^{\lambda +1}}K_\lambda (mx)=\frac{(x/m)^{\epsilon 1}}{(2\pi )^{2\epsilon }}K_{1\epsilon }(mx).$$
(191)
In this case we obtain
$$B_N(0,0,2,2,2,1)=\frac{2^{22\epsilon }}{(\mathrm{\Gamma }(1+\epsilon ))^3\mathrm{\Gamma }(2\epsilon )}_0^{\mathrm{}}\left(K_\epsilon (\xi )\right)^3K_{1\epsilon }(\xi )\xi ^{2+2\epsilon }๐\xi .$$
(192)
Using Eq. (187) one has
$$K_{1\epsilon }(\xi )=K_1(\xi )\frac{\epsilon }{\xi }K_0(\xi )+O(\epsilon ^2)$$
(193)
and
$$\frac{2^{22\epsilon }\left(K_\epsilon (\xi )\right)^3K_{1\epsilon }(\xi )\xi ^{2+2\epsilon }}{(\mathrm{\Gamma }(1+\epsilon ))^3\mathrm{\Gamma }(2\epsilon )}=4\xi ^2\left(1+\left(1\frac{1}{\xi }+\mathrm{ln}(e^{\gamma _E}\xi /2)\right)\epsilon +O(\epsilon ^2)\right).$$
(194)
The result reads
$`B_N(0,0,2,2,2,1)`$ $`=`$ $`4(1+\epsilon ){\displaystyle _0^{\mathrm{}}}\left(K_0(\xi )\right)^3K_1(\xi )\xi ^2๐\xi 4\epsilon {\displaystyle _0^{\mathrm{}}}\left(K_0(\xi )\right)^4\xi ๐\xi `$ (195)
$`+8\epsilon {\displaystyle _0^{\mathrm{}}}\left(K_0(\xi )\right)^3K_1(\xi )\xi ^2\mathrm{ln}(e^{\gamma _E}\xi /2)๐\xi .`$
This result is not yet written in terms of $`L_4(r)`$ and $`L_4^l(r)`$ and will be dealt with after introduction of the reduction procedure (The result for $`\epsilon =0`$ was given earlier in Sec. 3.6).
### 5.5 The example $`B_N(0,0,2,3,3,4)`$
In order to demonstrate the power of the configuration space technique also in a more complex setting we finally choose the diagram in Fig. 14(c) as an example. The modified propagators now read
$$D^{(2)}(x,m)=\frac{(x/m)^{1+\epsilon }}{8(2\pi )^{2\epsilon }}K_{1\epsilon }(mx),D^{(3)}(x,m)=\frac{(x/m)^{2+\epsilon }}{48(2\pi )^{2\epsilon }}K_{2\epsilon }(mx)$$
(196)
Using the expansions
$$K_{1\epsilon }(\xi )=K_1(\xi )+\frac{\epsilon }{\xi }K_0(\xi )+O(\epsilon ^2),K_{2\epsilon }(\xi )=K_2(\xi )+\frac{2\epsilon }{\xi }K_1(\xi )+\frac{2\epsilon }{\xi ^2}K_0(\xi )+O(\epsilon ^2)$$
(197)
we obtain
$`B_N(0,0,2,3,3,4)={\displaystyle \frac{2^{62\epsilon }}{3(\mathrm{\Gamma }(1+\epsilon ))^2\mathrm{\Gamma }(2\epsilon )}}{\displaystyle _0^{\mathrm{}}}K_\epsilon (\xi )\left(K_{1\epsilon }(\xi )\right)^2K_{2\epsilon }(\xi )\xi ^{7+2\epsilon }๐\xi `$ (198)
$`=`$ $`{\displaystyle \frac{1+\epsilon }{192}}{\displaystyle _0^{\mathrm{}}}K_0(\xi )\left(K_1(\xi )\right)^2K_2(\xi )\xi ^7๐\xi +{\displaystyle \frac{\epsilon }{96}}{\displaystyle _0^{\mathrm{}}}K_0(\xi )\left(K_1(\xi )\right)^2K_2(\xi )\xi ^6๐\xi `$
$`+{\displaystyle \frac{\epsilon }{96}}{\displaystyle _0^{\mathrm{}}}K_0(\xi )\left(K_1(\xi )\right)^3\xi ^6๐\xi +{\displaystyle \frac{\epsilon }{96}}{\displaystyle _0^{\mathrm{}}}\left(K_0(\xi )\right)^2\left(K_1(\xi )\right)^2\xi ^5๐\xi `$
$`+{\displaystyle \frac{\epsilon }{96}}{\displaystyle _0^{\mathrm{}}}K_0(\xi )\left(K_1(\xi )\right)^2K_2(\xi )\xi ^7\mathrm{ln}(e^{\gamma _E}\xi /2)๐\xi .`$
### 5.6 The reduction procedure
Especially in the last expression given by Eq. (198) there are many different integrals which differ from the basis $`L_4(r)`$ and $`L_4^l(r)`$. However, the integrands can be reduced to integrands involving $`K_0(\xi )`$ and $`K_1(\xi )`$ only by using the relation
$$K_n(\xi )=2\frac{n1}{\xi }K_{n1}(\xi )+K_{n2}(\xi ).$$
(199)
After the first step, namely the expansion of Bessel functions for non-integer indices, the above relation establishes the second step in our reduction procedure. Finally, we use
$`{\displaystyle \frac{d}{d\xi }}K_0(\xi )`$ $`=`$ $`K_1(\xi )\text{and}`$
$`{\displaystyle \frac{d}{d\xi }}K_1(\xi )`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(K_0(\xi )+K_2(\xi )\right)=K_0(\xi ){\displaystyle \frac{1}{\xi }}K_1(\xi )`$ (200)
to perform the third and last step. For instance one has
$`L_4^{(1)}(r)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}\left(K_0(\xi )\right)^3K_1(\xi )\xi ^r๐\xi ={\displaystyle _0^{\mathrm{}}}\left(K_0(\xi )\right)^3{\displaystyle \frac{dK_0(\xi )}{d\xi }}\xi ^r๐\xi =`$ (201)
$`=`$ $`\left[K_0(\xi )\left(K_0(\xi )\right)^3\xi ^r\right]_0^{\mathrm{}}+{\displaystyle _0^{\mathrm{}}}K_0(\xi ){\displaystyle \frac{d}{d\xi }}\left(K_0(\xi )\right)^3\xi ^r๐\xi =`$
$`=`$ $`3{\displaystyle _0^{\mathrm{}}}K_0(\xi )\left(K_0(\xi )\right)^2{\displaystyle \frac{dK_0(\xi )}{d\xi }}\xi ^r๐\xi +r{\displaystyle _0^{\mathrm{}}}K_0(\xi )\left(K_0(\xi )\right)^3\xi ^{r1}๐\xi =`$
$`=`$ $`3{\displaystyle _0^{\mathrm{}}}\left(K_0(\xi )\right)^3K_1(\xi )\xi ^r๐\xi +r{\displaystyle _0^{\mathrm{}}}\left(K_0(\xi )\right)^4\xi ^{r1}๐\xi `$
and therefore
$$L_4^{(1)}(r)=_0^{\mathrm{}}\left(K_0(\xi )\right)^3K_1(\xi )\xi ^r๐\xi =\frac{r}{4}_0^{\mathrm{}}\left(K_0(\xi )\right)^4\xi ^{r1}๐\xi =\frac{r}{4}L_4(r1).$$
(202)
The reduction formulas are given in general by
$`L_n^{(m)}(r)`$ $`=`$ $`{\displaystyle \frac{1}{nm+1}}\left((rm+1)L_n^{(m1)}(r1)(m1)L_n^{(m2)}(r)\right),`$ (203)
$`L_n^{l(m)}(r)`$ $`=`$ $`{\displaystyle \frac{1}{nm+1}}\left((rm+1)L_n^{l(m1)}(r1)+L_n^{(m1)}(r1)(m1)L_n^{l(m2)}(r)\right)`$
and are coded in MATHEMATICA in order to automatically reduce to the master integrals . After executing the second and third step of the recursion the results for the above examples read
$`B_N(0,0,2,2,2,2)`$ $`=`$ $`2(1+\epsilon )L_4(3)+4\epsilon L_4^l(3)+O(\epsilon ^2),`$
$`B_N(0,0,2,2,2,1)`$ $`=`$ $`2L_4(1)+4\epsilon L_4^l(1)+O(\epsilon ^2),`$
$`B_N(0,0,2,3,3,4)`$ $`=`$ $`{\displaystyle \frac{1}{36}}L_4(3){\displaystyle \frac{1}{144}}L_4(5){\displaystyle \frac{1}{576}}L_4(7)+O(\epsilon ).`$ (204)
Let us concentrate on the $`O(\epsilon ^0)`$ terms. By matching our numerical results to the results obtained by using the RECURSOR package ,
$`B_N(0,0,2,2,2,2)`$ $`=`$ $`{\displaystyle \frac{3}{8}}+{\displaystyle \frac{7}{16}}\zeta (3)+O(\epsilon ),`$
$`B_N(0,0,2,2,2,1)`$ $`=`$ $`{\displaystyle \frac{7}{4}}\zeta (3)+O(\epsilon ),`$
$`B_N(0,0,2,3,3,4)`$ $`=`$ $`{\displaystyle \frac{1}{576}}+O(\epsilon )`$ (205)
we can determine the master integrals $`L_4(r)`$ for a few values of its argument. One has
$$L_4(3)=\frac{3}{16}+\frac{7}{32}\zeta (3),L_4(1)=\frac{7}{8}\zeta (3),16L_4(3)2L_4(5)L_4(7)=1.$$
(206)
In particular the last relation is interesting because it implies a sum rule for the weighted product of three Bessel functions, namely
$$_0^{\mathrm{}}K_0(\xi )\left(K_1(\xi )\right)^2K_2(\xi )\xi ^7๐\xi =\frac{1}{3}.$$
(207)
This surprising identity has been cross-checked by numerical integration. Note that for the $`O(\epsilon ^0)`$ contribution only odd values of $`r`$ appear as arguments in the master integral $`L_4(r)`$.
### 5.7 Laplace-type asymptotics
The two master integrals Eq. (184) are the basic set for the reduction procedure. A first estimate of the numerical magnitude of these integrals can be easily inferred from the asymptotic expansion of the integrals at large $`q`$,
$`L_4(q)`$ $`=`$ $`{\displaystyle \frac{\pi ^2\mathrm{\Gamma }(2q)}{4^{2q+1}}}\left(1{\displaystyle \frac{1}{q1/2}}+O(1/q^2)\right),`$
$`L_4^l(q)`$ $`=`$ $`{\displaystyle \frac{\pi ^2\mathrm{\Gamma }(2q)}{4^{2q+1}}}\left(\mathrm{\Psi }(2q)+\gamma _E3\mathrm{ln}2\right)\left(1{\displaystyle \frac{1}{q1/2}}+O(1/q^2)\right)`$ (208)
where $`\mathrm{\Psi }(x)=\mathrm{\Gamma }^{}(x)/\mathrm{\Gamma }(x)`$ is the logarithmic derivative of the $`\mathrm{\Gamma }`$-function. The direct configuration space representation is the most convenient for numerical evaluation. To see this, we modify the last term of the integrals $`L_4(q)`$ and $`L_4^l(q)`$ by introducing parameters $`\kappa `$ and $`\kappa _l`$, respectively. In doing this, we find that the relations
$$L_4(q)=\frac{\pi ^2\mathrm{\Gamma }(2q)}{4^{2q+1}}\left(1\frac{1}{q+\kappa }\right),$$
(209)
$$L_4^l(q)=\frac{\pi ^2\mathrm{\Gamma }(2q)}{4^{2q+1}}\left(\mathrm{\Psi }(2q)+\gamma _E3\mathrm{ln}2\right)\left(1\frac{1}{q+\kappa _l}\right)$$
(210)
with $`\kappa =0.97`$ ($`\kappa _l=1.17`$) gives good results for the basis rational (resp. log-type) integrals with an accuracy better than $`1\%`$ for all values of the arguments $`q>1`$ ($`q>3`$). For $`q>3`$ ($`q>5`$) the relative accuracy is better than $`10^3`$. The leading order asymptotic formulas for the basis log-type integrals are less precise for small $`q`$ than its analogue in the rational case because the integrand in Eq. (184) is not positive for log-type integrals. Obviously, an exact solution via recurrence relations in momentum space will be much more complicated than these simple asymptotic formulae. Because the cancellation of significant figures in the numerical evaluation of terms with different transcendental structure is a dominating feature of the momentum space recurrence procedure for large $`q`$, the use of asymptotic formulas is worth considering.
### 5.8 Irreducible numerator: three-loop vacuum bubble case
Loop integrals may have numerator factors which involve the loop momenta and cannot be expanded in terms of the denominator pole factors. In such a case one speaks of irreducible numerator factors. Momentum space techniques can run into problems when non-trivial numerator factors appear. Take for example the numerator factor $`(k_1k_2)`$ for a $`n`$-loop bubble with $`n+1`$ massive lines. For $`n<3`$ the momentum space integral can be reduced to scalar integrals. For $`n=3`$, however, the numerator factorintegral is no longer reducible . The problem of irreducible numerator factors has a straightforward solution in configuration space by the integration-by-parts technique . The starting expression involving the non-trivial numerator factor $`(k_1k_2)`$ (or any other scalar product of linear independent inner moments) corresponds to
$$\stackrel{~}{\mathrm{\Pi }}_3^{}(0)=D(x,m)^2\left(_\mu D(x,m)\right)\left(^\mu (D(x,m))\right)d^Dx$$
(211)
where the asterix denotes the fact that we are dealing with a non-scalar master integral. Careful use of the integration-by-parts identities
$$_\mu \left(D(x,m)\mathrm{}D(x,m)\right)d^Dx=0,\text{ }\text{ }\text{ }\left(D(x,m)\mathrm{}D(x,m)\right)d^Dx=0$$
(212)
leads to the result
$$D(x,m)^2\left(_\mu D(x,m)\right)\left(^\mu (D(x,m))\right)d^Dx=\frac{1}{3}D(x,m)^3\text{ }\text{ }\text{ }D(x,m)d^Dx.$$
(213)
For the last integral we have used $`(\text{ }\text{ }\text{ }+m^2)D(x,m)=\delta (x)`$ and end up with
$$D(x,m)^3\text{ }\text{ }\text{ }D(x,m)d^Dx=m^2D(x,m)^4d^DxD(0,m)^3.$$
(214)
The value of $`D(0,m)`$ can be found for instance by integration in momentum space,
$$D(0,m)=\frac{1}{(2\pi )^D}\frac{d^Dp}{m^2+p^2}=\frac{2\pi ^{D/2}}{(2\pi )^D\mathrm{\Gamma }(D/2)}_0^{\mathrm{}}\frac{p^{D1}dp}{p^2+m^2}=\frac{m^{D2}}{(4\pi )^{D/2}}\mathrm{\Gamma }(1D/2).$$
(215)
In the end we obtain
$$\stackrel{~}{\mathrm{\Pi }}_3^{}(0)=\frac{m^2}{3}\stackrel{~}{\mathrm{\Pi }}_3(0)+\frac{1}{3}D(0,m)^3$$
(216)
where $`\stackrel{~}{\mathrm{\Pi }}_3(0)`$ is the scalar three-loop sunrise-type bubble diagram without a numerator factor.
### 5.9 Irreducible numerator: four-loop vacuum bubble case
Next we consider a four-loop diagram which appears as a second independent master integral $`V_2`$ of the sunrise topology in the classification of Ref. . In momentum space the second master integral $`V_2`$ has the additional numerator factor $`(k_1k_4)^2`$ as compared to the first scalar master integral $`V_1`$. The second master integral reads
$$\stackrel{~}{\mathrm{\Pi }}_4^{}(0)=\frac{(k_1k_4)^2(2\pi )^{4D}d^Dk_1d^Dk_2d^Dk_3d^Dk_4}{(m_1^2+k_1^2)(m_2^2+(k_2k_1)^2)(m_3^2+(k_3k_2)^2)(m_4^2+(k_4k_3)^2)(m_5^2+k_4^2)}.$$
(217)
Turning again to the equal mass case and using the configuration space representation, this integral can be written as
$$\stackrel{~}{\mathrm{\Pi }}_4^{}(0)=D(x,m)^3(_\mu _\nu D(x,m))(_\mu _\nu D(x,m))d^Dx.$$
(218)
It is apparent that by using integration-by-parts techniques this integral cannot be reduced to scalar integrals and/or integrals containing dโAlembertians. The easiest way to evaluate such an integral is to compute the derivatives directly. This is done with the help of the relation
$$\frac{1}{z}\frac{d}{dz}\left(z^\nu K_\nu (z)\right)=\left(z^{\nu 1}K_{\nu +1}(z)\right)$$
(219)
(cf. Eq. (A12)). Eq.(219) can be iterated and gives results for arbitrary high order derivatives of Bessel functions $`K_\lambda (z)`$ in terms of the same class of Bessel functions with shifted indices and powers in $`z`$. For the first derivative we obtain
$$_\mu D(x,m)=x_\mu \frac{m^{2\lambda +2}}{(2\pi )^{\lambda +1}}\frac{K_{\lambda +1}(mx)}{(mx)^{\lambda +1}}.$$
(220)
Since the resulting analytical expression for a given line of the diagram lies in the same class as the original line, the procedure of evaluating the integral is similar to the usual one. However, we cannot use the second derivative
$$_\mu _\nu D(x,m)=\frac{m^{2\lambda +2}}{(2\pi )^{\lambda +1}(mx)^{\lambda +1}}\left(g_{\mu \nu }K_{\lambda +2}(mx)\frac{x_\mu x_\nu }{x^2}K_{\lambda +1}(mx)\right)$$
(221)
directly under the integration sign. The reason is that the propagator has to be regarded as a distribution. There is a $`\delta `$-function singularity at the origin which is not taken into account in Eq. (221). Indeed, contracting the indices $`\mu `$ an $`\nu `$ in Eq. (221) one obtains
$$_\mu ^\mu D(x,m)=m^2D(x,m)$$
(222)
while the correct equation for the propagator reads $`(^2+m^2)D(x,m)=\delta (x)`$. Thus, a straightforward evaluation of derivatives is valid only for $`x0`$. The behaviour at the origin ($`x=0`$) requires special consideration. In practice, to treat this case one should not use higher order derivatives but stay at the level of the first derivative.
In order to deal with this case we introduce another master integral
$$\stackrel{~}{\mathrm{\Pi }}_4^{}(0)=D(x,m)_\mu D(x,m)_\nu D(x,m)^\mu D(x,m)^\nu D(x,m)d^Dx.$$
(223)
The relation between the two master integrals $`\mathrm{\Pi }_4^{}(0)`$ and $`\mathrm{\Pi }_4^{}(0)`$ is found to be
$$\stackrel{~}{\mathrm{\Pi }}_4^{}(0)=3\stackrel{~}{\mathrm{\Pi }}_4^{}(0)\frac{1}{8}m^4\stackrel{~}{\mathrm{\Pi }}_4(0)\frac{7}{8}m^2D(0,m)^4.$$
(224)
The quantity $`\stackrel{~}{\mathrm{\Pi }}_4^{}(0)`$ can be calculated with the explicit use of first order derivatives within our technique. The analytical result for the pole part reads
$$๐ฉ_4\stackrel{~}{\mathrm{\Pi }}_4^{}(0)=m^{10}\left(\frac{3}{8\epsilon ^4}\frac{277}{144\epsilon ^3}\frac{37837}{6912\epsilon ^2}\frac{4936643}{414720\epsilon }+O(\epsilon ^0)\right)$$
(225)
($`๐ฉ_4=((4\pi )^{2\epsilon }/\mathrm{\Gamma }(1=\epsilon ))^4`$, cf. Eq.(50)) and the $`\epsilon `$-expansion in the form
$`๐ฉ_4\stackrel{~}{\mathrm{\Pi }}_4^{}(0)`$ $`=`$ $`m^{10}(0.375\epsilon ^41.923611\epsilon ^35.474103\epsilon ^211.90356\epsilon ^1`$ (226)
$`27.99303104.5384\epsilon 663.6123\epsilon ^23703.241\epsilon ^3+O(\epsilon ^4)).`$
Since the results for $`\stackrel{~}{\mathrm{\Pi }}_4(0)`$, $`\stackrel{~}{\mathrm{\Pi }}_4^{}(0)`$, and $`D(0,m)^4`$ are known, Eq. (224) can be used to obtain the final result for the original integral,
$`๐ฉ_4\stackrel{~}{\mathrm{\Pi }}_4^{}(0)`$ $`=`$ $`m^{10}(1.6875\epsilon ^47.8125\epsilon ^321.20964\epsilon ^244.76955\epsilon ^1`$ (227)
$`97.07652290.9234\epsilon 1719.809\epsilon ^28934.731\epsilon ^3+O(\epsilon ^4))`$
which again verifies the result given in Ref. .
Differentiation of the massive propagator leads to expressions of a similar functional form which makes the configuration space technique a universal tool for calculating any master integral of the sunrise topology. This technique is also useful for finding master integrals. Indeed, new master integrals appear when there is a possibility to add new derivatives into the integrands which cannot be eventually removed by using the equations of motion or integration-by-parts recurrence relations. But once again: without explicit inclusion of the $`\delta `$-function only one derivative is allowed. Otherwise one misses tadpole parts of the result. Therefore, the new master integral should contain just one derivative for each line except for one line. This allows one to enumerate the number of non-scalar master integrals, i.e. master integrals including non-trivial numerator factors. For instance, in the five-loop case (six propagators) there will be only two non-scalar master integrals.
## 6 Generalization: Still within reach
Apart from the diagrams with sunrise-type topology treated in the previous sections there are more involved topologies close to the sunset-type topology that still allow for quite a simple numerical evaluation. We conclude this report by discussing several such cases which are still in reach of the methods presented here.
### 6.1 Generalization to the spectacle topology
In this subsection we give a formula for a more general topology which occurs when one propagator is removed from an initial three-loop, tetrahedron-shaped bubble diagram. In the original classification of Ref. these diagrams are class $`E`$ diagrams belonging to the spectacle topology (see Fig. 15). The formulas obtained in this subsection are efficient for the numerical integration of diagrams of the spectacle topology. We have, however, not been able to find an analytical solution to the problem. The main obstacle of generalizing the configuration space technique to a general multi-loop diagram is the angular integration. The configuration space technique proved to be rather successful for general diagrams in the massless case but it brings no essential simplification in the general massive case (see e.g. Ref. ). However, for special configurations the angular integration can be explicitly performed with a reasonably simple integrand left for the radial integration. Diagrams of the spectacle topology such as the ones shown in Fig. 15 are examples of such a configuration.
The configuration space expression of a spectacle topology diagram can be written in a form suitable for our purpose (see Fig. 15(a)),
$$D(xy,m)D(x,m)^2d^DxD(y,m)^2d^Dy.$$
(228)
The key relation which leads to the simplification of the configuration space integral with spectacle topology is an addition theorem for Bessel functions in Eq. (A17), allowing to perform the integration over the relative angle in the propagator $`D(xy,m)`$. Using Eq. (A17) for $`Z=H^+`$ and substituting $`m=e^{i\pi /2}`$, we can perform an analytic continuation in order to obtain a relation for the modified Bessel functions $`K`$ and $`I`$,
$`{\displaystyle \frac{K_\lambda (R)}{R^\lambda }}`$ $`=`$ $`2^\lambda \mathrm{\Gamma }(\lambda ){\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}(\lambda +k){\displaystyle \frac{I_{\lambda +k}(\rho )}{\rho ^\lambda }}{\displaystyle \frac{K_{\lambda +k}(r)}{r^\lambda }}C_k^\lambda (\mathrm{cos}\phi )`$ (229)
where we have $`R=\sqrt{r^2+\rho ^22r\rho \mathrm{cos}\phi }`$,
$`K_\lambda (z)`$ $`=`$ $`{\displaystyle \frac{i\pi }{2}}e^{\pi \lambda i/2}H_\lambda ^+(iz),`$
$`I_\lambda (z)`$ $`=`$ $`e^{\pi \lambda i/2}J_\lambda (e^{\pi i/2}z)\text{for }\pi <\text{arg}z{\displaystyle \frac{\pi }{2}},`$
$`I_\lambda (z)`$ $`=`$ $`e^{3\pi \lambda i/2}J_\lambda (e^{3\pi i/2}z)\text{for }{\displaystyle \frac{\pi }{2}}<\text{arg}z\pi .`$ (230)
Using the orthogonality relations for Gegenbauer polynomials in Eq. (B3), the sum disappears after integration over the relative angle and only one term contributes. We obtain
$`{\displaystyle \frac{K_\lambda (R)}{R^\lambda }๐\mathrm{\Omega }_\rho }`$ $`=`$ $`2^\lambda \mathrm{\Gamma }(\lambda ){\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}(\lambda +k){\displaystyle \frac{I_{\lambda +k}(\rho )}{\rho ^\lambda }}{\displaystyle \frac{K_{\lambda +k}(r)}{r^\lambda }}{\displaystyle C_k^\lambda (\mathrm{cos}\phi )๐\mathrm{\Omega }_\phi }`$ (231)
$`=`$ $`2^\lambda \mathrm{\Gamma }(\lambda )\lambda {\displaystyle \frac{I_\lambda (\rho )}{\rho ^\lambda }}{\displaystyle \frac{K_\lambda (r)}{r^\lambda }}{\displaystyle \frac{2\pi ^{\lambda +1}}{\mathrm{\Gamma }(\lambda +1)}}C_0^\lambda (1)`$
$`=`$ $`(2\pi )^{\lambda +1}{\displaystyle \frac{I_\lambda (\rho )}{\rho ^\lambda }}{\displaystyle \frac{K_\lambda (r)}{r^\lambda }},r>\rho ,`$
where the first equality is a consequence of the orthogonality relation with the trivial factor $`C_0^\lambda (1)=1`$. This result allows one to write down an expression for any spectacle-type diagram in the form of a two-fold integral with a simple integration measure,
$`{\displaystyle _0^{\mathrm{}}}D(x,m)^2x^{2\lambda +1}dx{\displaystyle _0^{\mathrm{}}}D(y,m)^2y^{2\lambda +1}dy\times `$ (232)
$`\left({\displaystyle \frac{K_\lambda (x)}{x^\lambda }}{\displaystyle \frac{I_\lambda (y)}{y^\lambda }}\theta (xy)+{\displaystyle \frac{K_\lambda (y)}{y^\lambda }}{\displaystyle \frac{I_\lambda (x)}{x^\lambda }}\theta (yx)\right)`$
$`=`$ $`{\displaystyle _0^{\mathrm{}}}D(x,m)^2I_\lambda (x)x^{\lambda +1}๐x{\displaystyle _x^{\mathrm{}}}D(y,m)^2K_\lambda (y)y^{\lambda +1}๐y`$
where $`\theta (x)`$ is the standard step-function distribution .
Note that the integration measure $`D(x,m)^2x^{2\lambda +1}dx`$ allows one to perform the integration by using efficient integration routines for numerical evaluation. The form of the weight function is close to $`e^{ax}x^\alpha `$ which suggests the use of Laguerre polynomials
$$L_n(x)=\frac{e^x}{n!}\frac{d^n}{dx^n}\left(x^ne^x\right)$$
(233)
as a convenient choice within the Gaussian numerical integration method. In fact, any modified propagator (with any power of the denominator) can be used as a factor in the integration measure $`D(x,m)^2x^{2\lambda +1}dx`$. This fact makes the representation universal and useful for the case of higher powers of denominators of the lines associated with pairs $`(x,0)`$ and $`(y,0)`$ of space-time points. If the angular structure of the diagram is preserved, the generalization to higher loops in the expressions for the radial measures is straightforward.
### 6.2 Example for the spectacle diagram
To show how this technique works, we give an example of the evaluation of a spectacle diagram. Consider an integer dimension space-time which, without loss of generality, can be chosen to be two-dimensional (an odd number of dimensions is trivial because, as we have seen earlier, the propagators degenerate to simple exponentials). The spectacle-type three-loop diagram can be obtained in a closed form. Indeed, in momentum space representation we have
$$S(M)=\frac{\stackrel{~}{\mathrm{\Pi }}(p)^2}{p^2+M^2}d^2p$$
(234)
for the basic spectacle diagram $`S`$ with $`\stackrel{~}{\mathrm{\Pi }}(p)`$ and the mass $`M`$ of the connecting propagator kept different. After the substitutions $`p=2m\mathrm{sinh}(\xi /2)`$ and $`t=e^\xi `$ we obtain
$$S(M)=\frac{1}{2\pi m^4}_0^1\frac{t\mathrm{ln}^2tdt}{(1t^2)(12\gamma t+t^2)}$$
(235)
where $`\gamma =1M^2/2m^2`$. Performing the integration, we finally obtain
$$S(M)=\frac{f(t_1)f(t_2)}{t_1t_2}$$
(236)
with $`t_{1,2}=\gamma \pm \sqrt{\gamma ^21}`$ and
$$f(t)=\frac{8t\mathrm{Li}_3(1/t)(t+7)\zeta (3)}{8\pi m^4(t^21)}.$$
(237)
$`\mathrm{Li}_3(z)`$ is the trilogarithm function
$$\mathrm{Li}_3(z)=\underset{k=1}{\overset{\mathrm{}}{}}\frac{z^k}{k^3},|z|<1.$$
(238)
For $`M=2m`$ (i.e. $`\gamma =1`$) the integral in Eq. (235) simplifies and one finds a simple answer in terms of the (for the present context) standard transcendental numbers $`\mathrm{ln}2`$ and $`\zeta (3)`$,
$$S(2m)=\frac{1}{4\pi m^4}\left(\frac{7}{8}\zeta (3)\mathrm{ln}2\right).$$
(239)
For the equal mass case $`M=m`$ we obtain a result in terms of Clausenโs trilogarithm $`\mathrm{Cl}_3(2\pi /3)`$. As one can conclude from these results, the conjugate pair of the sixth order roots of unity, $`\mathrm{exp}(\pm 2\pi i/3)`$ play an important role in this case again in accordance with the general analysis in . The origin of the appearance of the sixth order roots of unity as the parameters of the analytical expressions of the diagrams lies in the mismatch of masses along the lines of the diagrams. However, the exceptional case $`M=2m`$, where one line has twice the mass of the other lines (which results in the drastic simplification), also keeps us within the set of the sixth order roots of unity. The key parameter in this case is simply the natural number $`1`$ which is definitely one of the sixth order roots of unity.
Turning to the configuration space representation in Eq. (232) we find
$`S(M)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}xK_0(mx)^2dx{\displaystyle _0^{\mathrm{}}}yK_0(my)^2dy\times `$ (240)
$`\left(K_0(Mx)I_0(My)\theta (xy)+K_0(My)I_0(Mx)\theta (yx)\right).`$
for the basic spectacle diagram. An explicit numerical integration of Eq. (240) shows coincidence with the analytical result in Eq. (236) which we have checked numerically for arbitrary values of $`M`$. In this example the analytical result has a rather simple form. This is not true if high powers of the denominators enter. Then the corresponding one-loop insertions are rather cumbersome and an explicit integration in configuration space is more convenient.
### 6.3 About the occurrence of Clausenโs dilogarithms
We want to have a closer look at how the square of Clausenโs dilogarithm $`\mathrm{Cl}_2(\pi /3)^2`$ can emerge at the level of spectacle topology diagrams. As transcendental number, Clausenโs dilogarithm $`\mathrm{Cl}_2(\pi /3)`$ characterizes the analytical results for three-loop bubbles. Its presence was discovered in the impressive treatise of David Broadhurst on the role of the sixth order roots of unity for the transcendental structure of results for Feynman diagrams in quantum field theory .
We consider the spectacle diagram in the form shown in Fig. 15(b). Here the expression for the generalized middle line is a product of the one-loop propagator and the standard particle propagator. We express the one-loop propagator using a dispersion representation with the spectral density $`\rho (s)`$ and obtain
$$\frac{1}{p^2+m^2}\stackrel{~}{\mathrm{\Pi }}(p)=\frac{1}{p^2+m^2}_{4m^2}^{\mathrm{}}\frac{\rho (s)ds}{s+p^2}=_{4m^2}^{\mathrm{}}\frac{\rho (s)ds}{sm^2}\left(\frac{1}{p^2+m^2}\frac{1}{s+p^2}\right).$$
(241)
We take only the first term which is sufficient for obtaining the result we are aiming for. In this case the integral becomes independent of $`p^2`$ and can be considered separately. One has
$$I=\frac{1}{p^2+m^2}_{4m^2}^{\mathrm{}}\frac{\rho (s)ds}{sm^2}$$
(242)
which leads to the sunrise diagram after the two other line shown in Fig. 15(b) have been added with a normalization factor given by the integral. One factor $`\mathrm{Cl}_2(\pi /3)`$ results from integrating the overall sunrise diagram which is composed of the propagator $`(p^2+m^2)^1`$ from Eq. (242). The two other lines of the diagram are shown in Fig. 15(b). The second factor $`\mathrm{Cl}_2(\pi /3)`$ has to be found in the normalization factor given by the integral in Eq. (242). Note that the very structure of this contribution โ the square of a number which first appeared at the lower loop level โ suggests a hint for its search. It should emerge as an iteration of a lower order contribution in accordance with the iterative structure of the R-operation which provides a general framework for the analysis of multi-loop diagrams. The following consideration confirms this conjecture. Consider the quantity
$$N=_{4m^2}^{\mathrm{}}\frac{\rho (s)ds}{sm^2}$$
(243)
and take $`\rho (s)`$ to be the spectral density in $`D`$-dimensional space-time ,
$$\rho (s)=\frac{(s4m^2)^{\lambda 1/2}}{2^{4\lambda +1}\pi ^{\lambda +1/2}\mathrm{\Gamma }(\lambda +1/2)\sqrt{s}},\sqrt{s}>2m.$$
(244)
Next consider the first order contribution of the expansion in $`\epsilon `$ near the space-time dimension $`D_0=2`$. The expansion in $`\lambda =\epsilon `$ results in
$$\frac{(s4m^2)^{\epsilon 1/2}}{\mu ^{2\epsilon }\sqrt{s}}=\frac{1}{\sqrt{s(s4m^2)}}\left(1+\epsilon \mathrm{ln}\left(\frac{s4m^2}{\mu ^2}\right)+O(\epsilon ^2)\right).$$
(245)
Therefore, the relevant first order term in $`\epsilon `$ reads
$$\mathrm{\Delta }_\epsilon \rho (s)=\frac{\mathrm{ln}((s4m^2)/m^2)}{2\pi \sqrt{s(s4m^2)}}$$
(246)
where $`\mu =m`$ has been chosen for convenience. Next we change variables according to
$$\sqrt{s}=2m\mathrm{cosh}(\xi /2),t=e^\xi $$
(247)
to obtain
$$\mathrm{\Delta }_\epsilon \rho (4m^2\mathrm{cosh}^2(\xi /2))=\frac{(\mathrm{ln}t+2\mathrm{ln}(1t))t}{2\pi m^2(1t^2)}.$$
(248)
For the quantity in question we find
$$\mathrm{\Delta }_\epsilon N=_{4m^2}^{\mathrm{}}\frac{\mathrm{\Delta }_\epsilon \rho (s)ds}{sm^2}=\frac{1}{2\pi m^2}_0^1\frac{(\mathrm{ln}t+2\mathrm{ln}(1t))dt}{1+t+t^2}.$$
(249)
The roots of the denominator of the integrand in Eq. (249) are now given by $`t_{3,4}=\mathrm{exp}(\pm 2\pi i/3)`$ which again is a conjugate pair of the sixth order roots of unity. After integrating this equation we readily find
$`\mathrm{\Delta }_\epsilon N`$ $`=`$ $`{\displaystyle \frac{1}{2\pi m^2\sqrt{3}}}\left({\displaystyle \frac{\pi }{3}}\mathrm{ln}3+\mathrm{Im}\left(\mathrm{Li}_2\left(e^{2i\pi /3}\right)\mathrm{Li}_2\left(e^{2i\pi /3}\right)\right)\right)`$ (250)
$`=`$ $`{\displaystyle \frac{1}{\pi m^2\sqrt{3}}}\left(\mathrm{Cl}_2\left({\displaystyle \frac{2\pi }{3}}\right){\displaystyle \frac{\pi }{6}}\mathrm{ln}3\right).`$
Using the relation
$$\mathrm{Cl}_2\left(\frac{2\pi }{3}\right)=\frac{2}{3}\mathrm{Cl}_2\left(\frac{\pi }{3}\right)$$
(251)
we finally obtain
$$\mathrm{\Delta }_\epsilon N=\frac{2}{3\pi m^2\sqrt{3}}\left(\mathrm{Cl}_2\left(\frac{\pi }{3}\right)\frac{\pi }{4}\mathrm{ln}3\right).$$
(252)
Therefore, in the first order of the $`\epsilon `$ expansion of the spectacle diagram we indeed found the remarkable contribution proportional to $`\mathrm{Cl}_2(\pi /3)^2`$. In our calculation it emerges naturally as the iteration of the lower order term . Originally, the presence of this contribution has been conjectured and confirmed in by direct numerical computation of the finite part of the general three-loop bubble in four-dimensional space-time.
### 6.4 The insertion of massless irreducible loops
One can consider a sunrise-type diagram with one or more irreducible loops by which we mean a generalized line more complicated than an ordinary standard propagator or a (large) power of them. Consider, for instance, the replacement of a line by a subdiagram of fish-type topology as shown in Fig. 16. The configuration space technique leads to a numerical solution in this case because we can replace one of the propagator factors $`D(x,m)`$ by the two-point propagator in the irreducible subdiagram, using the fact that
$$\mathrm{\Pi }(x)=D(x,\sqrt{s})\rho (s)๐s$$
(253)
where $`\rho (s)`$ is the spectral density of the subdiagram. If the subdiagram is a massless fish-type diagram, the spectral density is given by
$$\rho (s)=6\zeta (3)\delta (s).$$
(254)
The known result in momentum space can be restored by using the dispersion relation (20). On the other hand, we obtain
$$D(x)=6\zeta (3)D(x,\sqrt{s})\delta (s)๐s=6\zeta (3)D(x,0)$$
(255)
which up to a factor is the usual massless propagator.
### 6.5 The insertion of massive irreducible loops
Insertions of massive irreducible loops were used analytically for the calculation of corrections to baryon correlators and numerically for the mixing of heavy neutral mesons .
The massive insertion can be incorporated mainly because the two-loop diagrams are well analyzed with any mass configurations . For the semi-massive fish-type one-loop diagram with one massive and one massless line we can use
$$\rho (s)=\frac{4}{(4\pi )^4s}\left(\mathrm{Li}_2\left(\frac{m^2}{s}\right)+\frac{1}{2}\mathrm{ln}\left(1\frac{m^2}{s}\right)\mathrm{ln}\left(\frac{m^2}{s}\right)\right).$$
(256)
But even in the case with two different masses a result for the spectral density is available,
$$\rho (s)=\frac{4}{(4\pi )^4s}\left(\mathrm{Li}_2(x_1)+\mathrm{Li}_2(x_2)\frac{1}{2}\mathrm{ln}(x_1)\mathrm{ln}(1x_1)\frac{1}{2}\mathrm{ln}(x_2)\mathrm{ln}(1x_2)\right)$$
(257)
where
$`x_1`$ $`=`$ $`{\displaystyle \frac{2m_1^2}{m_1^2+m_2^2s+\sqrt{(m_1^2+m_2^2s)^24m_1^2m_2^2}}},`$
$`x_2`$ $`=`$ $`{\displaystyle \frac{2m_2^2}{m_1^2+m_2^2s+\sqrt{(m_1^2+m_2^2s)^24m_1^2m_2^2}}}`$ (258)
## 7 Summary and conclusion
To conclude, we have presented a review of configuration space techniques for the calculation of sunrise-type diagrams. We have shown that the singular or pole parts of any sunrise-type diagram are calculable analytically in the simplest possible manner. For the finite parts we obtained one-dimensional integrals of well-known functions which is very convenient for numerical evaluation. We have presented many sample calculations of multi-loop sunrise-type diagrams and have compared them with momentum space results in the literature when available. We have found agreement in every case. The agreement provides for a mutual check of the results which have been derived using very different methods. In the asymptotic analysis have dealed with different expansions. We can conclude that any kind of expansions is given as an expansion in parameters of the propagators and thus leads to the expansion coefficients in a straightforward way. Finally, extensions to non-standard propagators and other exotic settings figured out to be feasible as well as some considerations for slightly different topologies. All in all, we can stress again that the benefit of configuration space techiques is the fact that there is โalmostโ no integration in sunrise-type diagrams, as Eq. (6) shows.
### Acknowledgements
We thank Kostja Chetyrkin for discussion, Robert Delbourgo for kind attention and enthusiasm in advertising $`x`$-space, Andrey Davydychev for communication and help in finding references, Andrey Grozin for providing us with the source code of RECURSOR, with which a part of the calculation was cross-checked, David Broadhurst for criticism and friendly remarks, and Giampiero Passarino for communication. SG thanks the Mainz xloops-GiNaC-group for permanent interest, valuable comments and numerical cross-checks. AAP thanks V.A. Matveev for encouragement, attention and support, and P. Baikov for illuminating discussions of the present status of the optimization of recurrence relations under study by him. This work was supported by the INTAS grant No. 03-51-4007. SG acknowledges support by the DFG as a guest scientist in Mainz, by the Estonian target financed project No. 0182647s04, and by the grant No. 6216 given by the Estonian Science Foundation. The work of AAP was supported in part by the Russian Fund for Basic Research under contract No. 03-02-17177 and by the grant NS-2184.2003.2.
## Appendix A Properties of Bessel functions
Since Bessel functions play a crucial role in our calculations we collect a number of definitions and properties of Bessel functions, mostly taken from the handbooks of Watson (1944) , Prudnikov, Brychkov and Marichev (1990) , and Gradshteyn and Ryshik (1994) .
### Bessel functions and Hankel functions
Ordinary Bessel functions are solutions of the Bessel differential equation
$$z^2\frac{d^2Z_\nu (z)}{dz^2}+z\frac{dZ_\nu (z)}{dz}+(z^2\nu ^2)Z_\nu (z)=0$$
(A1)
($`\nu `$ need not be an integer). There are two classes of solutions $`Z_\nu (z)`$. Bessel functions of the first kind $`J_\nu (z)`$ are nonsingular at the origin $`z=0`$ while Bessel functions of the second kind $`Y_\nu (z)`$ are singular at $`z=0`$. Combining Bessel functions of both classes, we end up with Hankel functions
$$H_\nu ^\pm (z)=J_\nu (z)\pm iY_\nu (z)$$
(A2)
(instead of $`H_\nu ^\pm (z)`$ the notation $`H_\nu ^{(1)}=H_\nu ^+`$ and $`H_\nu ^{(2)}=H_\nu ^{}`$ are quite common).
Bessel functions of the first kind have the series expansion
$$J_\nu (z)=\left(\frac{z}{2}\right)^\nu \underset{k=0}{\overset{\mathrm{}}{}}\frac{(z^2/4)^k}{k!\mathrm{\Gamma }(\nu +k+1)}=\frac{(z/2)^\nu }{\mathrm{\Gamma }(1+\nu )}\left(1\frac{(z/2)^2}{1+\nu }+\mathrm{}\right).$$
(A3)
Including a factor occuring in practical applications, we can also write
$$\left(\frac{z}{2}\right)^\nu J_\nu (z)=\underset{k=0}{\overset{\mathrm{}}{}}\frac{(z^2/4)^k}{k!\mathrm{\Gamma }(\nu +k+1)}=\frac{1}{\mathrm{\Gamma }(1+\nu )}\left(1\frac{(z/2)^2}{1+\nu }+\mathrm{}\right).$$
(A4)
$`\mathrm{\Gamma }(z)`$ is Eulerโs Gamma function. Functional equations are of help in order to relate Bessel functions of different degree. They are given by
$$zZ_{\nu 1}(z)+zZ_{\nu +1}(z)=2\nu Z_\nu (z),Z_{\nu 1}(z)Z_{\nu +1}(z)=2\frac{d}{dz}Z_\nu (z)$$
(A5)
where $`Z`$ is any of the Bessel functions $`J`$, $`Y`$, or $`H^\pm `$. As a consequence we obtain
$$\frac{d^k}{dz^k}\left(z^\nu J_\nu (z)\right)=z^\nu J_{\nu +k}(z).$$
(A6)
### Modified Bessel functions
We can continue Bessel functions into the complex plane, ending up with modified Bessel functions of the first kind $`I_\nu (z)`$ and modified Bessel functions of the second kind $`K_\nu (z)`$, sometimes also known as McDonald functions. $`I_\nu (z)`$ and $`K_\nu (z)`$ are solutions of the modified Bessel differential equation
$$z^2\frac{d^2Z_\nu (z)}{dz^2}+z\frac{dZ_\nu (z)}{dz}(z^2\nu ^2)Z_\nu (z)=0$$
(A7)
The modified Bessel function of the second kind
$$K_\nu (z)=\frac{\pi }{2}\frac{I_\nu (z)I_\nu (z)}{\mathrm{sin}(\pi \nu )},\mathrm{\Gamma }(\nu )\mathrm{\Gamma }(1\nu )=\frac{\pi }{\mathrm{sin}(\pi \nu )}$$
(A8)
can be expressed by the modified Bessel function of the first kind $`I_\nu (z)`$ with series expansion
$$I_\nu (z)=\left(\frac{z}{2}\right)^\nu \underset{k=0}{\overset{\mathrm{}}{}}\frac{(z^2/4)^k}{k!\mathrm{\Gamma }(\nu +k+1)}$$
(A9)
Therefore, one has
$$\left(\frac{z}{2}\right)^\nu K_\nu (z)=\frac{\mathrm{\Gamma }(\nu )}{2}\left[1+\frac{1}{1\nu }\left(\frac{z}{2}\right)^2\frac{\mathrm{\Gamma }(1\nu )}{\mathrm{\Gamma }(1+\nu )}\left(\frac{z}{2}\right)^{2\nu }\right]+O(z^4,z^{2+2\nu }).$$
(A10)
For Eq. (A5) one can we derive functional equations for the modified Bessel functions,
$`zI_{\nu 1}(z)zI_{\nu +1}(z)=2\nu I_\nu (z),`$ $`I_{\nu 1}(z)+I_{\nu +1}(z)=2{\displaystyle \frac{d}{dz}}I_\nu (z),`$
$`zK_{\nu 1}(z)zK_{\nu +1}(z)=2\nu K_\nu (z),`$ $`K_{\nu 1}(z)+K_{\nu +1}(z)=2{\displaystyle \frac{d}{dz}}K_\nu (z),`$ (A11)
For the differentiation of the modified Bessel function including an appropriate factor we derive
$$\frac{d}{dz}\left(z^\nu K_\nu (z)\right)=z^\nu K_{\nu +1}(z).$$
(A12)
Bessel functions can also be differentiated with respect to their index, if this index is given by an integer. For $`K_\nu (x)`$ we for instance obtain
$$\left[\frac{K_\nu (z)}{\nu }\right]_{\nu =\pm n}=\pm \frac{1}{2}n!\underset{k=0}{\overset{n1}{}}\left(\frac{z}{2}\right)^{kn}\frac{K_k(z)}{k!(nk)},n\{0,1,\mathrm{}\}.$$
(A13)
Finally, the modified Bessel function $`K_\nu `$ and the Hankel function $`H_\nu ^+`$ are related by
$$K_\nu (z)=\frac{\pi i}{2}e^{i\nu \pi /2}H_\nu ^+(iz).$$
(A14)
### Asymptotic behaviour
Asymptotic expansions determine the behaviour of the modified Bessel functions of the second kind and the Hankel functions $`H_\nu ^{}`$ when the arguments become large. For the two functions just related to each other we obtain
$$K_{\nu ,N}^{as}(z)=\left(\frac{\pi }{2z}\right)^{1/2}e^z\left[\underset{n=0}{\overset{N1}{}}\frac{(\nu ,n)}{(2z)^n}+\theta \frac{(\nu ,N)}{(2z)^N}\right],(\nu ,n):=\frac{\mathrm{\Gamma }(\nu +n1/2)}{n!\mathrm{\Gamma }(\nu n1/2)}$$
(A15)
and
$$H_{\nu ,N}^{as}(iz)=\left(\frac{2}{\pi z}\right)^{1/2}e^{z+i\nu \pi /2}\left[\underset{n=0}{\overset{N1}{}}\frac{(1)^n(\nu ,n)}{(2z)^n}+\theta \frac{(1)^N(\nu ,N)}{(2z)^N}\right]$$
(A16)
where $`\theta [0,1]`$ is chosen appropriately. For the asymptotic behaviour of $`I_\nu (z)`$ see Eq. (C16).
### Addition theorem
The addition theorem for Bessel functions is given by
$$\frac{Z_\nu (mR)}{R^\nu }=2^\nu m^\nu \mathrm{\Gamma }(\nu )\underset{k=0}{\overset{\mathrm{}}{}}(\nu +k)\frac{J_{\nu +k}(m\rho )}{\rho ^\nu }\frac{Z_{\nu +k}(mr)}{r^\nu }C_k^\nu (\mathrm{cos}\phi )$$
(A17)
where $`C_k^\nu `$ are the Gegenbauer polynomials (cf. Appendix B), $`Z`$ is any of the Bessel functions $`J`$, $`Y`$, or $`H^\pm `$, $`R=\sqrt{r^2+\rho ^22r\rho \mathrm{cos}\phi }`$, and $`r>\rho `$. For $`r<\rho `$ the arguments of the Bessel functions on the right hand side of Eq. (A17) should be interchanged.
## Appendix B Properties of Gegenbauer polynomials
Gegenbauer polynomials can be generated using the characteristic polynomial
$$(t^22tx+1)^\nu =\underset{n=0}{\overset{\mathrm{}}{}}t^nC_n^\nu (x).$$
(B1)
In particular, we have $`C_0^\nu (x)=1`$, $`C_1^\nu (x)=2\nu x`$ and
$$(n+1)C_{n+1}^\nu (x)=2(n+\nu )xC_n^\nu (x)(n+2\nu 1)C_{n1}^\nu (x).$$
(B2)
Gegenbauer polynomials obey the orthogonality relations
$$C_m^\nu (\widehat{x}_1\widehat{x}_2)C_n^\nu (\widehat{x}_2\widehat{x}_3)๐\mathrm{\Omega }_2=\frac{2\pi ^{\nu +1}}{\mathrm{\Gamma }(\nu +1)}\frac{\nu \delta _{mn}}{n+\nu }C_n^\nu (\widehat{x}_1\widehat{x}_3),๐\mathrm{\Omega }_2=\frac{2\pi ^{\nu +1}}{\mathrm{\Gamma }(\nu +1)}$$
(B3)
where $`\widehat{x}_i`$ are unit four-vectors and $`d\mathrm{\Omega }_i`$ is the angular part of the integration measure. Finally, we have
$$C_n^\nu (1)=\frac{\mathrm{\Gamma }(n+2\nu )}{n!\mathrm{\Gamma }(2\nu )}.$$
(B4)
## Appendix C Cuts and discontinuities
Functions which are continuous at least on parts of the real axis can be continued to the whole complex plane. However, if singularities occur on the real axis or elsewhere in the complex plane, this continuation need no longer be unique. Examples which occur in the context of our calculations are logarithms and polylogarithms but also powers with non-integer exponent and Bessel functions. We will deal with these special cases in this appendix.
### The case of logarithms
For $`x=0`$ the function $`\mathrm{ln}(x)`$ has a non-removable singularity while for negative real values it is not defined. However, we can continue to the negative real axis by following a path which circumvents the origin. It makes a difference whether we turn around the origin in the mathematically positive or negative sense. If we for instance start at the real value $`x`$ and turn around in the positive sense to continue to a value at $`x`$, we can parametrize this by writing $`z=xe^{i\phi }`$ where $`\phi `$ starts at $`\phi =0`$ and increases until it reaches the negative real axis for $`\phi =\pi `$. By doing so, we obtain
$$\mathrm{ln}(xe^{+i\pi })=\mathrm{ln}x+i\pi .$$
(C1)
However, we can reach the negative real axis as well by starting with $`\phi `$ again but going to negative values, reaching the negative real axis for $`\phi =\pi `$. In this case we obtain
$$\mathrm{ln}(xe^{i\pi })=\mathrm{ln}xi\pi .$$
(C2)
Obviously, the value โ$`\mathrm{ln}(x)`$โ is not unique. Indeed, we can even think of the continuation to the complex plane as a sheet which winds up and rises each time we turn around the origin, the so-called Riemannian sheet. Nevertheless, for only one turn we can calculate the difference between the values. This quantity is known as discontinuity, and it can be calculated actually for every complex value $`z`$ except for singular points. The definition is given by
$$\mathrm{Disc}f(z)=f(ze^{+i0})f(ze^{i0}),e^{\pm i0}=\underset{ฯต0}{lim}e^{\pm iฯต},ฯต>0$$
(C3)
where $`z=0`$ is the singular point of the function $`f(z)`$. If we apply this definition to the case of a logarithm with negative argument, we obtain
$$\mathrm{Disc}\mathrm{ln}(x)=\mathrm{ln}(xe^{+i0})\mathrm{ln}(xe^{i0})=\mathrm{ln}(xe^{i\pi })\mathrm{ln}(xe^{+i\pi })=2\pi i\theta (x)$$
(C4)
where $`\theta (x)`$ is the step function ($`\theta (x)=1`$ for $`x>0`$, $`\theta (x)=0`$ for $`x<0`$). Note that we have been careful in replacing $`xe^{\pm i0}`$ by $`xe^{i\pi }`$. The discontinuity of powers of logarithms can be calculated as well, for instance
$`\mathrm{Disc}\mathrm{ln}^2(x)`$ $`=`$ $`\mathrm{ln}^2(xe^{i\pi })\mathrm{ln}^2(xe^{+i\pi })`$ (C5)
$`=`$ $`\left((\mathrm{ln}xi\pi )^2(\mathrm{ln}x+i\pi )^2\right)\theta (x)=4\pi i\mathrm{ln}x\theta (x).`$
For later convenience we divide the discontinuity by the factor $`2\pi i`$ to obtain
$$\frac{1}{2\pi i}\mathrm{Disc}\mathrm{ln}(x)=\theta (x),\frac{1}{2\pi i}\mathrm{Disc}\mathrm{ln}^2(x)=2\mathrm{ln}x\theta (x),\mathrm{}$$
(C6)
### The case of polylogarithms
Polylogarithms can be defined iteratively by
$$\mathrm{Li}_n(x)=_0^x\frac{\mathrm{Li}_{n1}(x^{})}{x^{}}๐x^{},\mathrm{Li}_1(x)=\mathrm{ln}(1x).$$
(C7)
For all the polylogarithms a discontinuity occurs for $`x>1`$. We obtain
$$\mathrm{Li}_1(xe^{\pm i0})=\mathrm{ln}(1xe^{\pm i0})=\mathrm{ln}(1+xe^{i\pi })=\left(\mathrm{ln}(x1)\pm i\pi \right)\theta (x1)$$
(C8)
where we have extracted a factor $`e^{i\pi }`$ from the argument, and therefore
$$\frac{1}{2\pi i}\mathrm{Disc}\mathrm{Li}_1(x)=\theta (x1).$$
(C9)
In order to calculate the discontinuity of $`\mathrm{Li}_2(x)`$, for $`x>1`$ we split the integral into two parts,
$`\mathrm{Li}_2(xe^{\pm i0})`$ $`=`$ $`{\displaystyle _0^1}{\displaystyle \frac{\mathrm{ln}(1x^{})}{x^{}}}๐x^{}{\displaystyle _1^x}{\displaystyle \frac{\mathrm{ln}(1x^{})}{x^{}}}๐x^{}`$ (C10)
$`=`$ $`\mathrm{Li}_2(1){\displaystyle _1^x}{\displaystyle \frac{\mathrm{ln}(x^{}1)}{x^{}}}๐x\pm \pi i{\displaystyle _1^x}{\displaystyle \frac{dx^{}}{x^{}}}.`$
The real parts cancel out when taking the difference needed for the determination of the discontinuity. Contrary to this the imaginary parts add up and one obtains
$`{\displaystyle \frac{1}{2\pi i}}\mathrm{Disc}\mathrm{Li}_2(x)`$ $`=`$ $`{\displaystyle \frac{1}{2\pi i}}{\displaystyle \frac{dx^{}}{x^{}}\mathrm{Disc}\mathrm{Li}_1(x^{})}={\displaystyle _1^x}{\displaystyle \frac{dx^{}}{x^{}}}=\mathrm{ln}x\theta (x1),\text{in the same way}`$
$`{\displaystyle \frac{1}{2\pi i}}\mathrm{Disc}\mathrm{Li}_3(x)`$ $`=`$ $`{\displaystyle \frac{1}{2\pi i}}{\displaystyle \frac{dx^{}}{x^{}}\mathrm{Disc}\mathrm{Li}_2(x^{})}={\displaystyle _1^x}{\displaystyle \frac{dx^{}}{x^{}}}\mathrm{ln}x^{}={\displaystyle \frac{1}{2}}\mathrm{ln}^2x\theta (x1),\text{in general}`$
$`{\displaystyle \frac{1}{2\pi i}}\mathrm{Disc}\mathrm{Li}_n(x)`$ $`=`$ $`{\displaystyle \frac{1}{n1}}\mathrm{ln}^{n1}x\theta (x1).`$ (C11)
### The case of non-integer powers
For powers with exponents close to integer values such as $`(bax)^n`$ the calculation of the discontinuity is closely related to the calculation of the discontinuity of the logarithm because
$$x^\epsilon =1+\epsilon \mathrm{ln}x+O(\epsilon ^2).$$
(C12)
Similarly, the discontinuity of $`(bax)^n`$ can be obtained in the more general case in closed form. In order to calculate the discontinuity of a power $`(bax)^n`$ with non-integer exponent $`n`$ (we include the minus sign for later convenience) we first have to ponder where the discontinuities can appear. Obviously, this is the case for $`bax<0`$. Taking $`a>0`$ and $`b0`$ which is the case for all applications in this report, we obtain $`x>b/a`$. Therefore, we obtain
$`\mathrm{Disc}(bax)^n=(bax^{+i0})^n(bax^{i0})^n=(b+axe^{i\pi })^n(b+axe^{+i\pi })^n`$ (C13)
$`=`$ $`\left(e^{+i\pi n}(axb)^ne^{i\pi n}(axb)^n\right)\theta (axb)=2i\mathrm{sin}(\pi n)(axb)^n\theta (axb).`$
Using
$$\mathrm{sin}(\pi n)=\frac{\pi }{\mathrm{\Gamma }(n)\mathrm{\Gamma }(1n)}=\mathrm{sin}(\pi n)$$
(C14)
where $`\mathrm{\Gamma }(x)`$ is Eulerโs Gamma function, we finally obtain
$$\frac{1}{2\pi i}\mathrm{Disc}(bax)^n=\frac{(axb)^n\theta (axb)}{\mathrm{\Gamma }(n)\mathrm{\Gamma }(1n)}$$
(C15)
Note that for negative integer values of $`n`$ the discontinuity vanishes.
### The case of Bessel functions
The final item about ambiguities found in the continuation into the complex plane we want to mention here is the so-called Stokesโ phenomenon which occurs for the asymptotic expansion of Bessel functions and is mentioned already in the handbook of Watson . The asymptotic expansion of the modified Bessel function of the first kind reads
$`I_\nu (z)`$ $`=`$ $`{\displaystyle \frac{e^z}{\sqrt{2\pi z}}}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^k}{(2z)^k}}{\displaystyle \frac{\mathrm{\Gamma }(\nu +k+1/2)}{\mathrm{\Gamma }(\nu k+1/2)k!}}`$ (C16)
$`+{\displaystyle \frac{e^{z\pm (\nu +1/2)\pi i}}{\sqrt{2\pi z}}}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(2z)^k}}{\displaystyle \frac{\mathrm{\Gamma }(\nu +k+1/2)}{\mathrm{\Gamma }(\nu k+1/2)k!}},`$
where the plus sign is valid for $`\pi /2<\mathrm{arg}z<3\pi /2`$ while the minus sign has to be taken for $`3\pi /2<\mathrm{arg}z<\pi /2`$. The surprising fact is that there are two different asymptotic expansions available at every point of the complex plane, except for the arguments $`\pm i\pi /2`$. Therefore, in continuing $`I_\nu (\pm ix)=I_\nu (xe^{\pm i\pi /2})`$ to the complex plane, the choice is still unique.
## Appendix D Singular contributions up to four loops
In the following we present complete results for the singular parts of sunrise-type diagrams with arbitrary masses up to four-loop order. The results are given in the $`\overline{\mathrm{MS}}`$-scheme in the Euclidean domain.
$`\stackrel{~}{\mathrm{\Pi }}_1^s(p,m_1,m_2)={\displaystyle \frac{1}{\epsilon }},`$
$`\stackrel{~}{\mathrm{\Pi }}_2^s(p,m_1,m_2,m_3)={\displaystyle \frac{1}{2\epsilon ^2}}{\displaystyle \underset{i}{}}m_i^2{\displaystyle \frac{1}{4\epsilon }}\left(p^2+2{\displaystyle \underset{i}{}}m_i^2(32\mathrm{}_i)\right),`$
$`\stackrel{~}{\mathrm{\Pi }}_3^s(p,m_1,m_2,m_3,m_4)={\displaystyle \frac{1}{6\epsilon ^3}}{\displaystyle \underset{ij}{}}m_i^2m_j^2+{\displaystyle \frac{1}{12\epsilon ^2}}\left(p^2{\displaystyle \underset{i}{}}m_i^2{\displaystyle \underset{i}{}}m_i^4+2{\displaystyle \underset{ij}{}}m_i^2m_j^2\left(43\mathrm{}_i\right)\right)+`$
$`+{\displaystyle \frac{1}{72\epsilon }}\left(2p^4+9p^2{\displaystyle \underset{i}{}}m_i^2(32\mathrm{}_i)9{\displaystyle \underset{i}{}}m_i^4(52\mathrm{}_i)+6{\displaystyle \underset{ij}{}}m_i^2m_j^2(2024\mathrm{}_i+3\mathrm{}_i^2+6\mathrm{}_i\mathrm{}_j)\right),`$
$`\stackrel{~}{\mathrm{\Pi }}_4^s(p,m_1,m_2,m_3,m_4,m_5)={\displaystyle \frac{1}{24\epsilon ^4}}{\displaystyle \underset{ijk}{}}m_i^2m_j^2m_k^2+`$ (D1)
$`{\displaystyle \frac{1}{48\epsilon ^3}}\left(p^2{\displaystyle \underset{ij}{}}m_i^2m_j^2{\displaystyle \underset{ij}{}}(m_i^4m_j^2+m_i^2m_j^4)+2{\displaystyle \underset{ijk}{}}m_i^2m_j^2m_k^2(54\mathrm{}_i)\right)`$
$`{\displaystyle \frac{1}{288\epsilon ^2}}(2p^4{\displaystyle \underset{i}{}}m_i^26p^2{\displaystyle \underset{i}{}}m_i^4+2{\displaystyle \underset{i}{}}m_i^6+3p^2{\displaystyle \underset{ij}{}}m_i^2m_j^2(118\mathrm{}_i)`$
$`6{\displaystyle \underset{ij}{}}(m_i^4m_j^2+m_i^2m_j^4)(74\mathrm{}_i)+12{\displaystyle \underset{ijk}{}}m_i^2m_j^2m_k^2(1520\mathrm{}_i+2\mathrm{}_i^2+6\mathrm{}_i\mathrm{}_j))`$
$`{\displaystyle \frac{1}{1728\epsilon }}(3p^6+2p^4{\displaystyle \underset{i}{}}m_i^2(3524\mathrm{}_i)18p^2{\displaystyle \underset{i}{}}m_i^4(218\mathrm{}_i)+2{\displaystyle \underset{i}{}}m_i^6(7724\mathrm{}_i)`$
$`+9p^2{\displaystyle \underset{ij}{}}m_i^2m_j^2(7188\mathrm{}_i+8\mathrm{}_i^2+24\mathrm{}_i\mathrm{}_j)216{\displaystyle \underset{ij}{}}(m_i^4m_j^2m_i^2m_j^4)\mathrm{}_i`$
$`18{\displaystyle \underset{ij}{}}(m_i^4m_j^2+m_i^2m_j^4)(4956\mathrm{}_i+4\mathrm{}_i^2+12\mathrm{}_i\mathrm{}_j)`$
$`+24{\displaystyle \underset{ijk}{}}m_i^2m_j^2m_k^2(105180\mathrm{}_i+30\mathrm{}_i^2+90\mathrm{}_i\mathrm{}_j2\mathrm{}_i^318\mathrm{}_i^2\mathrm{}_j12\mathrm{}_i\mathrm{}_j\mathrm{}_k))`$
where $`\mathrm{}_i=\mathrm{ln}(m_i^2/\mu ^2)`$. The indices $`i`$, $`j`$, and $`k`$ run over all mass indices. One can check that the general results listed in this Appendix reproduce the equal mass results in the main text.
## Appendix E Subtraction terms for the small $`x`$ singularities
The leading singularity at small $`x`$ is given by the massless propagator of the form
$$D(x,0)=\frac{\mathrm{\Gamma }(\lambda )}{4\pi ^{\lambda +1}x^{2\lambda }}.$$
(E1)
The next order of the small $`x`$-expansion for the propagator $`D(x,m)`$ is explicitly given by
$$D_1(x,0)=\frac{1}{4\pi ^{\lambda +1}x^{2\lambda }}\left(\left(\frac{x}{2}\right)^2\frac{\mathrm{\Gamma }(\lambda )}{1\lambda }\left(\frac{x}{2}\right)^{2\lambda }\frac{\mathrm{\Gamma }(1\lambda )}{\lambda }\right).$$
(E2)
This term is suppressed relative to the first term by one power of $`x^2`$ at small $`x`$ in four-dimensional space-time (however, this is not the case for two-dimensional space-time with $`\lambda =0`$). The term
$$D_2(x,0)=\frac{1}{4\pi ^{\lambda +1}x^{2\lambda }}\left(\frac{x}{2}\right)^2\left(\left(\frac{x}{2}\right)^2\frac{\mathrm{\Gamma }(\lambda )}{2(1\lambda )(2\lambda )}\left(\frac{x}{2}\right)^{2\lambda }\frac{\mathrm{\Gamma }(1\lambda )}{\lambda (\lambda +1)}\right)$$
(E3)
is further suppressed by one power of $`x^2`$ at small $`x`$. Therefore, the full subtraction of the three terms gives a rather smooth behaviour at small $`x`$ which is sufficient to obtain a regular integrand for the numerical integration.
## Appendix F Analytical results for the four-loop sunrise diagram
In this appendix we present some more details of our calculations for the four-loop sunrise diagram. For the analytical evaluation we take the last two terms of the integrand from Eq. (78). One has to integrate a product of two Bessel functions multiplied with powers of $`x`$ which can be done analytically. The explicit expression for the $`\epsilon `$-expansion of that part of the integral which is evaluated analytically reads
$`\stackrel{~}{\mathrm{\Pi }}_4^{\mathrm{ana}}(0)`$ $`=`$ $`m^6\{{\displaystyle \frac{5}{2\epsilon ^4}}{\displaystyle \frac{35}{3\epsilon ^3}}{\displaystyle \frac{4565}{144\epsilon ^2}}{\displaystyle \frac{58345}{864\epsilon }}`$ (F1)
$`{\displaystyle \frac{1456940638037}{7779240000}}{\displaystyle \frac{17099\pi ^2}{24}}{\displaystyle \frac{3857\pi ^4}{10}}+{\displaystyle \frac{2525968\zeta (3)}{105}}`$
$`+({\displaystyle \frac{55171475321621447}{1633640400000}}+{\displaystyle \frac{2457509\pi ^2}{144}}{\displaystyle \frac{1292537\pi ^4}{175}}`$
$`+{\displaystyle \frac{6752474831\zeta (3)}{44100}}+16530\pi ^2\zeta (3)+59508\zeta (5))\epsilon `$
$`+({\displaystyle \frac{10610679621089130529}{68612896800000}}+{\displaystyle \frac{92781949\pi ^2}{864}}{\displaystyle \frac{4290113759\pi ^4}{110250}}{\displaystyle \frac{22591\pi ^6}{14}}`$
$`+{\displaystyle \frac{952412727629\zeta (3)}{9261000}}+244476\pi ^2\zeta (3)168606\zeta (3)^2+{\displaystyle \frac{32210272\zeta (5)}{35}})\epsilon ^2`$
$`+({\displaystyle \frac{5963907632629558995931}{14408708328000000}}+{\displaystyle \frac{1325204033\pi ^2}{5184}}`$
$`{\displaystyle \frac{464379085699\pi ^4}{6615000}}{\displaystyle \frac{48529231\pi ^6}{2205}}`$
$`{\displaystyle \frac{312138383154103\zeta (3)}{277830000}}+{\displaystyle \frac{7285043\pi ^2\zeta (3)}{6}}+43529\pi ^4\zeta (3){\displaystyle \frac{238229084\zeta (3)^2}{105}}`$
$`+{\displaystyle \frac{13583011297\zeta (5)}{2940}}+247950\pi ^2\zeta (5)+1190160\zeta (7))\epsilon ^3+O(\epsilon )^4\}.`$
This expression shows the real complexity of the calculation which reveals itself in the structure of the results. The main feature is that the terms cannot be simultaneously simplified to all orders in $`\epsilon `$. By a special choice of the normalization factor one can make the leading term and, in fact, even all pole terms simple, but then the higher order terms contain rather lengthy combinations of transcendental numbers that are not reducible in terms of standard quantities such as the Riemann $`\zeta `$-functions. Note also that the rational coefficients of transcendental numbers are very big and there is a huge numerical cancellation between the rational and transcendental parts of the answer (see also the discussion in Ref. ).
## Appendix G Integrand for the numerical integration
For the numerical evaluation we take the first three terms of the integrand from Eq. (78). One has to integrate them numerically as there is a product of three or more Bessel functions which is too complicated to be done analytically. To find the $`\epsilon `$-expansion of the integral one has to first expand the integrand in $`\epsilon `$. The expression for the $`\epsilon `$-expansion is quite lengthy. We therefore give explicit results only for $`\epsilon =0`$. For this part the integrand for the numerical integration over $`z=mx`$ reads
$`\mathrm{\Pi }_4^{\mathrm{num}}(x)`$ $`=`$ $`m^6(66108l144l^2+192l^3+{\displaystyle \frac{384}{z^6}}{\displaystyle \frac{384}{z^4}}+{\displaystyle \frac{768l}{z^4}}+{\displaystyle \frac{24}{z^2}}{\displaystyle \frac{480l}{z^2}}+{\displaystyle \frac{576l^2}{z^2}}`$ (G1)
$`{\displaystyle \frac{111z^2}{16}}+{\displaystyle \frac{147lz^2}{2}}117l^2z^2+24l^3z^2+24l^4z^2+{\displaystyle \frac{165z^4}{32}}`$
$`+{\displaystyle \frac{201lz^4}{16}}+{\displaystyle \frac{9l^2z^4}{2}}24l^3z^4+12l^4z^4+{\displaystyle \frac{75z^6}{512}}{\displaystyle \frac{405lz^6}{128}}+{\displaystyle \frac{531l^2z^6}{64}}`$
$`{\displaystyle \frac{15l^3z^6}{2}}+{\displaystyle \frac{9l^4z^6}{4}}+{\displaystyle \frac{375z^8}{2048}}{\displaystyle \frac{825lz^8}{1024}}+{\displaystyle \frac{315l^2z^8}{256}}{\displaystyle \frac{51l^3z^8}{64}}+{\displaystyle \frac{3l^4z^8}{16}}`$
$`+{\displaystyle \frac{1875z^{10}}{131072}}{\displaystyle \frac{375lz^{10}}{8192}}+{\displaystyle \frac{225l^2z^{10}}{4096}}{\displaystyle \frac{15l^3z^{10}}{512}}+{\displaystyle \frac{3l^4z^{10}}{512}})K_1(z)`$
$`+m^6({\displaystyle \frac{512}{z^5}}+{\displaystyle \frac{384}{z^3}}{\displaystyle \frac{768l}{z^3}}+{\displaystyle \frac{24}{z}}+{\displaystyle \frac{288l}{z}}{\displaystyle \frac{384l^2}{z}}52z+120lz64l^3z`$
$`{\displaystyle \frac{15z^3}{8}}21lz^3+48l^2z^324l^3z^3+{\displaystyle \frac{75z^5}{32}}{\displaystyle \frac{135lz^5}{16}}+9l^2z^53l^3z^5`$
$`+{\displaystyle \frac{125z^7}{512}}{\displaystyle \frac{75lz^7}{128}}+{\displaystyle \frac{15l^2z^7}{32}}{\displaystyle \frac{l^3z^7}{8}})K_1(z)^2+m^6{\displaystyle \frac{128K_1(z)^5}{z^2}}.`$
Here $`l=\mathrm{ln}(e_E^\gamma z/2)`$, $`z=mx`$, and $`\gamma _E=\mathrm{\Gamma }^{}(1)=0.577\mathrm{}`$ is Eulerโs constant. As shown in Fig. 5, the plot of this function as well as the shapes of the corresponding functions in higher orders of $`\epsilon `$ are very smooth and quite similar. The analytical expressions for higher orders in $`\epsilon `$, however, become much longer. Note that the new functions $`f_n(z)`$ first appear at order $`\epsilon ^2`$.
The smoothness of the zeroth order integrand as shown in Eq. (G1) implies that the numerical integration is quite easy to execute. Because the integrand vanishes exponentially for large values of $`z`$ and has no singularities of the kind $`z\mathrm{ln}z`$ for small values of $`z`$, the integration can in principle range from $`0`$ to $`\mathrm{}`$. However, for practical reasons we had to instruct Wolframโs Mathematica system for symbolic manipulations (which we used for all of our calculations presented here) that the integrand vanishes for $`z=0`$. On the other hand, the asymptotic expansion of the integrand together with the integration measure is dominated by the part
$$\frac{6\pi ^2m^2}{512}z^{10}\mathrm{ln}^4(e^\gamma z/2)K_1(z)z^3dz(K_\lambda (z)\sqrt{\frac{\pi }{2z}}e^z\text{ for }z\mathrm{}).$$
(G2)
Integrated from $`\mathrm{\Lambda }`$ to $`\mathrm{}`$, this part gives a contribution
$$\frac{6\pi ^2m^2}{512}\mathrm{\Lambda }^{25/2}\mathrm{ln}^4(e^\gamma \mathrm{\Lambda }/2)e^\mathrm{\Lambda }$$
(G3)
and terms which are of subleading order. Therefore, the result can be well estimated by
$$\stackrel{~}{\mathrm{\Pi }}_4^{\mathrm{num}}(0)=2\pi ^2_0^{\mathrm{}}\mathrm{\Pi }_4^{\mathrm{num}}(x)x^3๐x2\pi ^2_0^\mathrm{\Lambda }\mathrm{\Pi }_4^{\mathrm{num}}(x)x^3๐x+\frac{6\pi ^2m^2}{512}\mathrm{\Lambda }^{25/2}\mathrm{ln}^4(e^\gamma \mathrm{\Lambda }/2)e^\mathrm{\Lambda }$$
(G4)
and $`\mathrm{\Lambda }`$ can be adjusted to any desired precision.
A possibility to avoid any cutoff is to change the integration variable such that the interval $`[0,\mathrm{}]`$ is mapped onto $`[0,1]`$. Then the integration can be done numerically with the additional information that the integrand vanishes identically at both end points. Possible transformations of this kind are for instance $`z=\mathrm{ln}(1/t)`$ or $`z=(1t)/t`$ for $`t[0,1]`$.
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# Quantum geons
## 1 Introduction
A number of models describing particles as solitonsโstable self-bound concentrations of field energyโhave been proposed. In an early paper Rosen obtained soliton solutions from the interaction of a complex scalar field and the electromagnetic field. Later Cooperstock and Rosen obtained soliton solutions by coupling a complex scalar field to both electromagnetism and gravity. In the meantime Wheeler , using an entirely different theoretical framework, obtained soliton solutions (which he dubbed geons) from a purely classical model of electromagnetism coupled to gravity. More recently Moroz, Penrose and Tod obtained soliton solutions of the Schrรถdinger-Newton equations. See and references therein for a window into the literature.
The model presented here is similar to that of Cooperstock and Rosen.<sup>2</sup><sup>2</sup>2Our Eqs. (22)โ(25) are equivalent to Eqs. (5.3)โ(5.6) of . What differs is the treatment of boundary conditions and the introduction of certain constraints. Cooperstock and Rosen require all fields to be finite and continuous at the center of the particle and they integrate the field equations outward. Parameters are adjusted until the asymptotic wave function vanishes. The resulting model contains several adjustable parameters.
In contrast, we impose asymptotic boundary conditions and integrate the field equations inward, toward the center of the particle. And we require the asymptotically measured mass and charge to equal the mass and charge parameters of the action. The resulting model contains no adjustable parameters.
In Section 2 we write down the action for a complex scalar field minimally coupled to the electromagnetic and gravitational fields, and use the principle of stationary action to derive the corresponding field equations (the Klein-Gordon, Einstein and Maxwell equations). Then we propose a stationary spherically symmetric trial solution for the metric tensor, wave function and vector potential and substitute it into the field equations. We obtain a system of four coupled nonlinear differential equations, as well as two integral constraints related to the conservation of charge and energy. In Section 3 we search for localized particle-like solutions. After calculating the asymptotic behavior of the equations and imposing asymptotic boundary conditions, we are left with a single adjustable parameterโthe system electric charge. By again appealing to the principle of stationary action we find that the charge must vanish. Hence, we arrive at an eigensystem with no adjustable parameters. Using numerical methods we calculate the eigenmodes and find five massive particle species without charge or spin. We name these particles quantum geons and show that the spacetime curvature diverges at the center of each geon. In Section 4 we explore the implications of this singularity and argue that it is benign. In Section 5 we discuss our results.
## 2 Field equations
Consider the action
$$S=\frac{1}{16\pi }_V\left[๐ฑ๐ฅ^{\mu \nu }๐ฅ_{\mu \nu }+\frac{1}{2}\left(\overline{\mathrm{\Psi }^{:\mu }}\mathrm{\Psi }_{:\mu }+\mathrm{\Psi }^{:\mu }\overline{\mathrm{\Psi }_{:\mu }}\right)M^2\overline{\mathrm{\Psi }}\mathrm{\Psi }\right]\sqrt{d^4x}$$
(1)
for a complex scalar field $`\mathrm{\Psi }`$ within a four-dimensional volume $`V`$ with spacetime curvature $`๐ฑ`$, where a colon denotes the generalized covariant derivative \[9, p. 167\]
$$\mathrm{\Psi }_{:\mu }\mathrm{\Psi }_{;\mu }+iQ๐ _\mu \mathrm{\Psi },$$
(2)
a semicolon denotes the covariant derivative of general relativity, an overline denotes complex conjugation, $`Q`$ is electric charge, $`๐ _\mu `$ is the electromagnetic vector potential,
$$๐ฅ_{\mu \nu }๐ _{\mu ;\nu }๐ _{\nu ;\mu }$$
(3)
is the electromagnetic field tensor and $`M`$ is mass (assumed positive). All quantities are expressed in natural units where the speed of light, Planckโs reduced constant ($`\mathrm{}`$), Newtonโs gravitational constant and Coulombโs electrostatic constant are unity. Our notation and conventions are detailed in Appendix A.
Make small arbitrary variations $`\delta S`$ in the action by making small arbitrary variations $`\delta ๐_{\alpha \beta }`$, $`\delta ๐ _\mu `$, $`\delta \mathrm{\Psi }`$ and $`\delta \overline{\mathrm{\Psi }}`$ in the fields, and neglect boundary terms arising from integration by parts. Then impose the condition $`\delta S=0`$ for stationary action. The terms proportional to $`\delta \mathrm{\Psi }`$ and $`\delta \overline{\mathrm{\Psi }}`$ yield the Klein-Gordon equation
$$\mathrm{\Psi }_{}^{:\mu }{}_{:\mu }{}^{}+M^2\mathrm{\Psi }=0$$
(4)
and its complex conjugate. The terms proportional to $`\delta ๐_{\alpha \beta }`$ yield the Einstein equations
$$๐ฆ^{\alpha \beta }=8\pi ๐ณ^{\alpha \beta },$$
(5)
where
$$๐ฆ^{\alpha \beta }๐ฑ^{\alpha \beta }\frac{1}{2}๐^{\alpha \beta }๐ฑ$$
(6)
is the Einstein tensor and
$`๐ณ^{\alpha \beta }`$ $``$ $`{\displaystyle \frac{1}{16\pi }}\left(\overline{\mathrm{\Psi }^{:\alpha }}\mathrm{\Psi }^{:\beta }+\mathrm{\Psi }^{:\alpha }\overline{\mathrm{\Psi }^{:\beta }}๐^{\alpha \beta }\overline{\mathrm{\Psi }^{:\mu }}\mathrm{\Psi }_{:\mu }+M^2๐^{\alpha \beta }\overline{\mathrm{\Psi }}\mathrm{\Psi }\right)`$ (7)
$`{\displaystyle \frac{1}{4\pi }}\left(๐ฅ_{}^{\alpha }{}_{\mu }{}^{}๐ฅ^{\beta \mu }\frac{1}{4}๐^{\alpha \beta }๐ฅ^{\mu \nu }๐ฅ_{\mu \nu }\right)`$
is the energy-momentum tensor. The terms proportional to $`\delta ๐ _\mu `$ yield the inhomogeneous Maxwell equations
$$๐ฅ_{}^{\mu \nu }{}_{;\nu }{}^{}=4\pi ๐ฉ^\mu ,$$
(8)
where
$$๐ฉ^\mu \frac{iQ}{16\pi }\left(\overline{\mathrm{\Psi }}\mathrm{\Psi }^{:\mu }\overline{\mathrm{\Psi }^{:\mu }}\mathrm{\Psi }\right)$$
(9)
is the electromagnetic current vector. The homogeneous Maxwell equations
$$๐ฅ_{\mu \nu ;\alpha }+๐ฅ_{\alpha \mu ;\nu }+๐ฅ_{\nu \alpha ;\mu }=0$$
(10)
follow immediately from the definition of $`๐ฅ_{\mu \nu }`$.
Let the volume $`V`$ in Eq. (1) include all of space over some arbitrary time interval $`T`$. We will later impose boundary conditions which require $`๐ฑ`$, $`๐ _\mu `$ and $`\mathrm{\Psi }`$ to vanish asymptotically. Thus we are justified in neglecting boundary terms when integrating by parts in the expression for $`\delta S`$.
When Eqs. (4), (5) and (8) are satisfied the action is stationary and the scalar curvature is
$$๐ฑ=8\pi ๐ณ=2M^2\overline{\mathrm{\Psi }}\mathrm{\Psi }\overline{\mathrm{\Psi }^{:\mu }}\mathrm{\Psi }_{:\mu },$$
(11)
where $`๐ณ๐ณ_\mu ^\mu `$. Substituting for $`๐ฑ`$ in Eq. (1) gives
$$S=\frac{1}{16\pi }_V\left(M^2\overline{\mathrm{\Psi }}\mathrm{\Psi }๐ฅ^{\mu \nu }๐ฅ_{\mu \nu }\right)\sqrt{d^4x}$$
(12)
as the stationary value of the action.
### 2.1 Trial solution
Our goal is to solve Eqs. (4), (5) and (8) for the fields $`๐_{\mu \nu }`$, $`\mathrm{\Psi }`$ and $`๐ _\mu `$. In order to make the calculations tractable we consider a static spherically symmetric spacetime with spherical polar coordinates $`x^0=t`$, $`x^1=r`$, $`x^2=\theta `$ and $`x^3=\varphi `$, and metric tensor
$$๐_{\mu \nu }=\left(\begin{array}{cccc}u& 0& 0& 0\\ 0& v& 0& 0\\ 0& 0& r^2& 0\\ 0& 0& 0& r^2\mathrm{sin}^2\theta \end{array}\right),$$
(13)
where $`u`$ and $`v`$ are real functions of $`r`$. The corresponding spherically symmetric wave function and vector potential are assumed to have the forms
$$\mathrm{\Psi }=Re^{i\omega t}$$
(14)
and
$$Q๐ _\mu =\left(\begin{array}{c}\mathrm{\Phi }\\ 0\\ 0\\ 0\end{array}\right),$$
(15)
where $`R`$ and $`\mathrm{\Phi }`$ are real functions of $`r`$ and $`\omega `$ is a real constant.
Expressions (13), (14) and (15) for $`๐_{\mu \nu }`$, $`\mathrm{\Psi }`$ and $`๐ _\mu `$ comprise a trial solution of field equations (4), (5) and (8) and allow us to calculate all quantities of interest in terms of $`R`$, $`u`$, $`v`$ and $`\mathrm{\Phi }`$. We will show that the field equations reduce to a set of four coupled nonlinear differential equations in $`R`$, $`u`$, $`v`$ and $`\mathrm{\Phi }`$. Then we will solve these equations subject to physically plausible boundary conditions and constraints. The fact that the field equations are, indeed, satisfied by the assumed trial solution is its ultimate justification.
### 2.2 Einstein, energy-momentum and electromagnetic tensors
Given trial solution (13), (14) and (15) we can calculate the quantities in field equations (4), (5) and (8). The nonzero elements of the Einstein and energy-momentum tensors are
$`๐ฆ_{0}^{}{}_{}{}^{0}`$ $`=`$ $`{\displaystyle \frac{v^{}}{rv^2}}{\displaystyle \frac{1}{r^2}}+{\displaystyle \frac{1}{r^2v}}`$
$`๐ฆ_{1}^{}{}_{}{}^{1}`$ $`=`$ $`{\displaystyle \frac{u^{}}{ruv}}{\displaystyle \frac{1}{r^2}}+{\displaystyle \frac{1}{r^2v}}`$
$`๐ฆ_{2}^{}{}_{}{}^{2}`$ $`=`$ $`{\displaystyle \frac{u^{\prime \prime }}{2uv}}{\displaystyle \frac{\left(u^{}\right)^2}{4u^2v}}{\displaystyle \frac{u^{}v^{}}{4uv^2}}+{\displaystyle \frac{u^{}}{2ruv}}{\displaystyle \frac{v^{}}{2rv^2}}`$
$`๐ฆ_{3}^{}{}_{}{}^{3}`$ $`=`$ $`๐ฆ_{2}^{}{}_{}{}^{2}`$ (16)
and
$`๐ณ_{0}^{}{}_{}{}^{0}`$ $`=`$ $`{\displaystyle \frac{M^2R^2}{16\pi }}+{\displaystyle \frac{\left(\omega \mathrm{\Phi }\right)^2R^2}{16\pi u}}+{\displaystyle \frac{\left(R^{}\right)^2}{16\pi v}}+{\displaystyle \frac{\left(\mathrm{\Phi }^{}\right)^2}{8\pi Q^2uv}}`$
$`๐ณ_{1}^{}{}_{}{}^{1}`$ $`=`$ $`{\displaystyle \frac{M^2R^2}{16\pi }}{\displaystyle \frac{\left(\omega \mathrm{\Phi }\right)^2R^2}{16\pi u}}{\displaystyle \frac{\left(R^{}\right)^2}{16\pi v}}+{\displaystyle \frac{\left(\mathrm{\Phi }^{}\right)^2}{8\pi Q^2uv}}`$
$`๐ณ_{2}^{}{}_{}{}^{2}`$ $`=`$ $`{\displaystyle \frac{M^2R^2}{16\pi }}{\displaystyle \frac{\left(\omega \mathrm{\Phi }\right)^2R^2}{16\pi u}}+{\displaystyle \frac{\left(R^{}\right)^2}{16\pi v}}{\displaystyle \frac{\left(\mathrm{\Phi }^{}\right)^2}{8\pi Q^2uv}}`$
$`๐ณ_{3}^{}{}_{}{}^{3}`$ $`=`$ $`๐ณ_{2}^{}{}_{}{}^{2},`$ (17)
where primes denote differentiation with respect to $`r`$. The contraction of the energy-momentum tensor is
$$๐ณ=\frac{M^2R^2}{4\pi }\frac{\left(\omega \mathrm{\Phi }\right)^2R^2}{8\pi u}+\frac{\left(R^{}\right)^2}{8\pi v}.$$
(18)
The nonzero elements of the electromagnetic field tensor are
$$Q๐ฅ_{01}=Q๐ฅ_{10}=\mathrm{\Phi }^{}$$
(19)
and the only nonzero element of the electromagnetic current vector is
$$๐ฉ^0=\frac{Q\left(\omega \mathrm{\Phi }\right)R^2}{8\pi u}.$$
(20)
We also note that
$$Q^2๐ฅ^{\mu \nu }๐ฅ_{\mu \nu }=\frac{2\left(\mathrm{\Phi }^{}\right)^2}{uv}.$$
(21)
### 2.3 Differential equations
The Klein-Gordon equation (4) becomes
$$R^{\prime \prime }+\left(\frac{2}{r}+\frac{u^{}}{2u}\frac{v^{}}{2v}\right)R^{}+\left[\frac{\left(\omega \mathrm{\Phi }\right)^2}{u}M^2\right]vR=0.$$
(22)
There are four Einstein equations (5), one for each nonzero component of $`๐ฆ_{\mu }^{}{}_{}{}^{\nu }`$. However, the equations corresponding to $`๐ฆ_{2}^{}{}_{}{}^{2}`$ and $`๐ฆ_{3}^{}{}_{}{}^{3}`$ are identical. Furthermore, they can be derived by combining the $`๐ฆ_{0}^{}{}_{}{}^{0}`$, $`๐ฆ_{1}^{}{}_{}{}^{1}`$ and Klein-Gordon equations. So we retain the $`๐ฆ_{0}^{}{}_{}{}^{0}`$, $`๐ฆ_{1}^{}{}_{}{}^{1}`$ and Klein-Gordon equations, and discard the $`๐ฆ_{2}^{}{}_{}{}^{2}`$ and $`๐ฆ_{3}^{}{}_{}{}^{3}`$ equations. It will prove convenient to work with the equations obtained by adding and subtracting the $`๐ฆ_{0}^{}{}_{}{}^{0}`$ and $`๐ฆ_{1}^{}{}_{}{}^{1}`$ equations:
$$\frac{u^{}}{2u}\frac{v^{}}{2v}+\frac{1v}{r}=rv\left[\frac{M^2R^2}{2}+\frac{\left(\mathrm{\Phi }^{}\right)^2}{Q^2uv}\right]$$
(23)
and
$$\frac{u^{}}{2u}+\frac{v^{}}{2v}=\frac{rv}{2}\left[\frac{\left(\omega \mathrm{\Phi }\right)^2R^2}{u}+\frac{\left(R^{}\right)^2}{v}\right].$$
(24)
The inhomogeneous Maxwell equations (8) yield the Poisson equation
$$\mathrm{\Phi }^{\prime \prime }+\left(\frac{2}{r}\frac{u^{}}{2u}\frac{v^{}}{2v}\right)\mathrm{\Phi }^{}+\frac{Q^2R^2v}{2}\left(\omega \mathrm{\Phi }\right)=0.$$
(25)
The four coupled nonlinear differential equations (22) through (25) are to be solved for the functions $`R`$, $`u`$, $`v`$ and $`\mathrm{\Phi }`$ subject to physical constraints and boundary conditions.
### 2.4 Conserved quantities
In this section we derive expressions for the electric charge, energy, angular momentum and stationary action of our model.
#### 2.4.1 Charge
The electromagnetic current vector satisfies the conservation law
$$๐ฉ_{}^{\mu }{}_{;\mu }{}^{}=\left(๐ฉ^\mu \right)_{,\mu }=0.$$
(26)
The corresponding conserved quantity
$$Q=k_q_{\mathrm{}}๐ฉ^0\sqrt{drd\theta d\varphi }$$
(27)
is the electric charge, where $`k_q`$ is a real constant (to be determined later). Substitute $`๐ฉ^0`$ from Eq. (20) and let
$$=r^2\sqrt{\left|uv\right|}\mathrm{sin}\theta $$
(28)
then integrate over $`\theta `$ and $`\varphi `$ to get
$$1=\frac{k_q}{2}_0^{\mathrm{}}\frac{r^2\left(\omega \mathrm{\Phi }\right)R^2}{u}\sqrt{|uv|}๐r.$$
(29)
Another expression for $`Q`$ can be obtained as follows. Substituting from Eq. (8) into Eq. (27) and using the identity
$$๐ฅ^{\mu \nu }{}_{;\nu }{}^{}=(๐ฅ^{\mu \nu })_{,\nu }$$
(30)
gives
$$Q=\frac{k_q}{4\pi }_{\mathrm{}}\left(๐ฅ^{0\nu }\right)_{,\nu }๐r๐\theta ๐\varphi .$$
(31)
For our static spherically symmetric metric the integrand is
$$\left(๐ฅ^{0\nu }\right)_{,\nu }=\frac{d}{dr}\left(\frac{\mathrm{\Phi }^{}}{Quv}\right),$$
(32)
so
$$Q^2=\frac{k_q}{4\pi }_{\mathrm{}}\frac{d}{dr}\left(\frac{\mathrm{\Phi }^{}}{uv}\right)๐r๐\theta ๐\varphi .$$
(33)
Substituting expression (28) for $``$ then integrating over $`\theta `$ and $`\varphi `$ gives
$`Q^2`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d}{dr}}\left({\displaystyle \frac{k_qr^2\mathrm{\Phi }^{}\sqrt{|uv|}}{uv}}\right)๐r`$ (34)
$`=`$ $`\stackrel{~}{Q}^2(\mathrm{}),`$
where
$$\stackrel{~}{Q}^2(r)\frac{k_qr^2\mathrm{\Phi }^{}\sqrt{|uv|}}{uv}$$
(35)
and we assume (to be justified later) $`\stackrel{~}{Q}^2(0)=0`$. The quantity $`\stackrel{~}{Q}^2(r_b)\stackrel{~}{Q}^2(r_a)`$ is the square of the charge contained in the region between $`r_a`$ and $`r_b`$.
#### 2.4.2 Energy
In our static spherically symmetric spacetime the vector field
$$\xi ^\mu =\left(\begin{array}{c}1\\ 0\\ 0\\ 0\end{array}\right)$$
(36)
satisfies the Killing equation
$$\xi ^{\mu ;\nu }+\xi ^{\nu ;\mu }=0.$$
(37)
If we let
$$๐ฉ_\xi ^\mu \xi _{}^{\mu ;\nu }{}_{;\nu }{}^{}$$
(38)
then the anti-symmetry of $`\xi ^{\mu ;\nu }`$ implies
$$๐ฉ_\xi ^\mu =\left(\xi ^{\mu ;\nu }\right)_{,\nu }$$
(39)
and
$$\left(๐ฉ_\xi ^\mu \right)_{,\mu }=0.$$
(40)
Thus $`๐ฉ_\xi ^\mu `$ is a conserved current associated with the time invariance of the metric. The corresponding conserved quantity
$$M=k_\xi _{\mathrm{}}๐ฉ_\xi ^0\sqrt{drd\theta d\varphi }$$
(41)
is the system mass (energy), where $`k_\xi `$ is a real constant (to be determined later). Calculation of $`๐ฉ_\xi ^\mu `$ from Eqs. (36) and (38) for our static spherically symmetric spacetime gives
$$๐ฉ_\xi ^\mu =\left(\begin{array}{c}๐ฑ_{0}^{}{}_{}{}^{0}\\ 0\\ 0\\ 0\end{array}\right).$$
(42)
But
$$๐ฑ_{0}^{}{}_{}{}^{0}=8\pi \left(๐ณ_{0}^{}{}_{}{}^{0}\frac{1}{2}๐ณ\right)$$
(43)
so Eq. (41) becomes
$$M=8\pi k_\xi _{\mathrm{}}\left(๐ณ_{0}^{}{}_{}{}^{0}\frac{1}{2}๐ณ\right)\sqrt{drd\theta d\varphi }.$$
(44)
Substituting expressions (17), (18) and (28) for $`๐ณ_{0}^{}{}_{}{}^{0}`$, $`๐ณ`$ and $``$ into (44), then integrating over $`\theta `$ and $`\varphi `$ gives
$$M=4\pi k_\xi _0^{\mathrm{}}r^2\left[\frac{\left(\omega \mathrm{\Phi }\right)^2R^2}{u}\frac{M^2R^2}{2}+\frac{\left(\mathrm{\Phi }^{}\right)^2}{Q^2uv}\right]\sqrt{|uv|}๐r.$$
(45)
Another expression for $`M`$ can be obtained as follows. Substituting from Eq. (39) into Eq. (41) gives
$$M=k_\xi _{\mathrm{}}\left(\xi ^{0;\nu }\right)_{,\nu }๐r๐\theta ๐\varphi .$$
(46)
For our static spherically symmetric metric the integrand is
$$\left(\xi ^{0;\nu }\right)_{,\nu }=\frac{d}{dr}\left(\frac{u^{}}{2uv}\right),$$
(47)
so
$$M=k_\xi _{\mathrm{}}\frac{d}{dr}\left(\frac{u^{}}{2uv}\right)๐r๐\theta ๐\varphi .$$
(48)
Substituting expression (28) for $``$ then integrating over $`\theta `$ and $`\varphi `$ gives
$`M`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d}{dr}}\left({\displaystyle \frac{2\pi k_\xi r^2u^{}\sqrt{|uv|}}{uv}}\right)๐r`$ (49)
$`=`$ $`\stackrel{~}{M}(\mathrm{}),`$
where
$$\stackrel{~}{M}(r)\frac{2\pi k_\xi r^2u^{}\sqrt{|uv|}}{uv}$$
(50)
and we assume (to be justified later) $`\stackrel{~}{M}(0)=0`$. The quantity $`\stackrel{~}{M}(r_b)\stackrel{~}{M}(r_a)`$ is the mass contained in the region between $`r_a`$ and $`r_b`$.
#### 2.4.3 Angular momentum
Based on the Killing field
$$\chi ^\mu =\left(\begin{array}{c}0\\ 0\\ 0\\ 1\end{array}\right)$$
(51)
and using an approach similar to that of Section 2.4.2 one can show that the system angular momentum is zero, as expected for a spherically symmetric model.
#### 2.4.4 Action
From Eqs. (12), (14), (21) and (28) the stationary value of the action is
$$S=\frac{T}{2}_0^{\mathrm{}}r^2\left[\frac{M^2R^2}{2}+\frac{\left(\mathrm{\Phi }^{}\right)^2}{Q^2uv}\right]\sqrt{|uv|}๐r.$$
(52)
## 3 Solution of the field equations
### 3.1 Solution at infinity
Our goal in this section is to determine the asymptotic ($`r\mathrm{}`$) expressions for the fields $`R`$, $`u`$, $`v`$ and $`\mathrm{\Phi }`$. We seek a particle-like solution with charge and energy distributions localized near the origin and we impose the following boundary conditions:
1. The wave function $`\mathrm{\Psi }`$ and all components of the vector potential $`๐ _\mu `$ asymptotically vanish.
2. Spacetime is asymptotically flat and local measurements by an asymptotic observer agree with the predictions of flat spacetime physics. Specifically, the asymptotic behaviors of $`u`$ and $`v`$ reproduce Newtonian gravity, and the asymptotic expressions for $`\mathrm{\Phi }`$ and $`R`$ satisfy the flat spacetime Poisson and Klein-Gordon equations.
These boundary conditions will guide us to candidate expressions for the asymptotic fields. We will then verify by direct substitution that these candidates asymptotically satisfy field equations (22) through (25).
The asymptotic expressions
$`u`$ $`=`$ $`12Mr^1+Q^2r^2`$
$`v`$ $`=`$ $`u^1`$ (53)
correspond to the (Reissner-Nordstrรถm) spacetime structure outside a charged nonrotating black hole. They yield flat spacetime (Newtonian) gravity as $`r\mathrm{}`$, as required by boundary condition 2. They also imply \[see Eq. (50)\]
$$\stackrel{~}{M}(\mathrm{})=4\pi k_\xi M.$$
(54)
Comparing this with Eq. (49) gives
$$k_\xi =\frac{1}{4\pi }.$$
(55)
The asymptotic expression
$$\mathrm{\Phi }=Q^2r^1$$
(56)
follows from boundary conditions 1 and 2 and the flat spacetime laws of electrostatics. Expressions (53) and (56) imply \[see Eq. (35)\]
$$\stackrel{~}{Q}^2(\mathrm{})=k_qQ^2.$$
(57)
Comparing this with Eq. (34) gives
$$k_q=1.$$
(58)
The asymptotic form of $`R`$ follows from boundary condition 2 and the asymptotic flat spacetime Klein-Gordon equation \[obtained by letting $`u,v1`$ in Eq. (22)\]
$$R^{\prime \prime }+\frac{2}{r}R^{}+\left[\left(\omega \mathrm{\Phi }\right)^2M^2\right]R=0,$$
(59)
where $`\mathrm{\Phi }`$ is given by (56). The solution which asymptotically vanishes (boundary condition 1) is
$$R=\frac{R_{\mathrm{}}e^{kr}}{r^{1+\sigma }},$$
(60)
where $`R_{\mathrm{}}`$ is a real constant,
$$\sigma =\frac{Q^2\omega }{k},$$
(61)
and
$$k=\sqrt{M^2\omega ^2}$$
(62)
with $`\omega ^2<M^2`$.
We now have asymptotic expressions for the fields which satisfy the boundary conditions. But we still need to verify that these expressions are consistent with our field equations. This is accomplished by substituting the asymptotic expressions for $`R`$, $`u`$, $`v`$ and $`\mathrm{\Phi }`$ into Eqs. (22) through (25) and verifying that the dominant terms satisfy the equations as $`r\mathrm{}`$. This procedure reveals that Eqs. (23) through (25) are, indeed, asymptotically satisfied. However, Eq. (22) is satisfied only if we let (see Appendix B)
$$\sigma =\frac{Q^2\omega }{k}+\frac{M}{k}\left(M^22\omega ^2\right),$$
(63)
which differs from (61). In order to satisfy boundary condition 2 we must make (61) and (63) consistent, which requires
$$\omega =\pm M/\sqrt{2}$$
(64)
and \[from (62)\]
$$k=M/\sqrt{2}.$$
(65)
In summary, based on the imposed boundary conditions we expect the asymptotic solution to have the form
$`R`$ $`=`$ $`R_{\mathrm{}}r^{1Q^2}e^{Mr/\sqrt{2}}`$
$`u`$ $`=`$ $`12Mr^1+Q^2r^2`$
$`v`$ $`=`$ $`u^1`$
$`\mathrm{\Phi }`$ $`=`$ $`Q^2r^1.`$ (66)
### 3.2 Physical constraints
In addition to satisfying field equations (22) through (25) the solution must also satisfy the charge and energy constraints of Section 2.4. Using results from Section 3.1 in Eqs. (29) and (45) we obtain the constraints
$$1=\frac{1}{2}_0^{\mathrm{}}r^2\left(\pm \frac{M}{\sqrt{2}}\mathrm{\Phi }\right)\frac{R^2}{u}\sqrt{|uv|}๐r$$
(67)
and
$$1=\frac{1}{M}_0^{\mathrm{}}r^2\left[(\pm \frac{M}{\sqrt{2}}\mathrm{\Phi })^2\frac{R^2}{u}\frac{M^2R^2}{2}+\frac{\left(\mathrm{\Phi }^{}\right)^2}{Q^2uv}\right]\sqrt{|uv|}๐r.$$
(68)
### 3.3 Solution strategy
If we are given values of the three parameters $`Q`$, $`M`$ and $`R_{\mathrm{}}`$ and a choice of sign in Eq. (64) then the asymptotic fields are fully defined by Eqs. (66). A complete solution can be obtained by integrating field equations (22) through (25) from the asymptotic limit ($`r\mathrm{}`$) down to the origin ($`r0`$). Since the two constraints of Section 3.2 must be satisfied, only one of the three parameters is independent. It will be convenient to take $`Q`$ as the independent parameter and regard $`M`$ and $`R_{\mathrm{}}`$ as functions of $`Q`$.
Once a solution is found for given $`Q`$ then the action $`S`$ can be calculated from Eq. (52). We may regard the action as a function of $`Q`$. The field equations, constraints and boundary conditions actually depend on $`Q^2`$, not $`Q`$, so we may write the action as $`S(Q^2)`$. Solutions of physical interest have stationary action so they must satisfy
$$\frac{S}{Q}=2QS^{}\left(Q^2\right)=0.$$
(69)
### 3.4 Solution for zero charge
We will consider solutions with $`Q=0`$. These solutions satisfy Eq. (69) so they are of physical interest.
With $`Q=0`$ the asymptotic fields (66) become
$`R`$ $`=`$ $`R_{\mathrm{}}r^1e^{Mr/\sqrt{2}}`$
$`u`$ $`=`$ $`12Mr^1`$
$`v`$ $`=`$ $`u^1`$
$`\mathrm{\Phi }`$ $`=`$ $`0.`$ (70)
Field equation (25) becomes
$$\mathrm{\Phi }^{\prime \prime }+\left(\frac{2}{r}\frac{u^{}}{2u}\frac{v^{}}{2v}\right)\mathrm{\Phi }^{}=0$$
(71)
which, given the asymptotic expression for $`\mathrm{\Phi }`$, yields the trivial solution (valid everywhere)
$$\mathrm{\Phi }=0.$$
(72)
Then field equations (22) through (24) become
$$R^{\prime \prime }+\left(\frac{2}{r}+\frac{u^{}}{2u}\frac{v^{}}{2v}\right)R^{}+\left(\frac{1}{2u}1\right)M^2vR=0,$$
(73)
$$\frac{u^{}}{2u}\frac{v^{}}{2v}+\frac{1v}{r}=\frac{M^2rvR^2}{2}$$
(74)
and
$$\frac{u^{}}{2u}+\frac{v^{}}{2v}=\frac{rv}{2}\left[\frac{M^2R^2}{2u}+\frac{\left(R^{}\right)^2}{v}\right],$$
(75)
and constraints (67) and (68) can be written
$$\mathrm{\Delta }2\sqrt{2}M_0^{\mathrm{}}\frac{r^2R^2}{u}\sqrt{|uv|}๐r=0$$
(76)
and
$$\pm 2\sqrt{2}2M_0^{\mathrm{}}r^2R^2\sqrt{|uv|}๐r=0.$$
(77)
Since both $`M`$ and the integral in (77) are positive this constraint can only be satisfied if we choose the upper sign. So we need only consider the upper signs in Eqs. (76) and (77), and \[from Eq. (64)\] $`\omega =M/\sqrt{2}`$.
Given values of $`M`$ and $`R_{\mathrm{}}`$, and starting with the asymptotic fields (70), Eqs. (73) through (75) can be numerically integrated all the way to the origin. The values of $`M`$ and $`R_{\mathrm{}}`$ can be adjusted until constraints (76) and (77) are satisfied.
Figure 1 shows the locus in the $`M`$$`R_{\mathrm{}}`$ plane which satisfies constraint (77). Figure 2 plots $`\mathrm{\Delta }`$ versus $`M`$ along that locus. Since $`\mathrm{\Delta }=0`$ when constraint (76) is satisfied the five zero crossings in Fig. 2 correspond to the solutions we seek. These solutions describe uncharged spinless massive particles. Because the masses appear spontaneously as eigenvalues of a self-gravitating quantum field (โmass without massโ) we borrow terminology from Wheeler and dub these particles quantum geons. We will index the five geons by the integer $`n=0,1,2,3,4`$ in order of increasing mass.
Various features of the geons are plotted in Figs. 3 through 5 and numerical characteristics are tabulated in Table 1. The probability density is
$$P\frac{\overline{\mathrm{\Psi }}\mathrm{\Psi }}{_{\mathrm{}}\overline{\mathrm{\Psi }}\mathrm{\Psi }\sqrt{drd\theta d\varphi }}=NR^2,$$
(78)
where the normalization factor
$$N=\frac{M}{8\pi \left(\sqrt{2}1\right)}$$
(79)
is determined from Eq. (77). The radial probability density
$$P_r=4\pi Nr^2R^2\sqrt{|uv|}$$
(80)
is plotted in Fig. 6. The stationary value of the action calculated from Eqs. (52) and (77) is
$$S=\frac{\left(\sqrt{2}1\right)TM}{2}$$
(81)
and the action per cycle ($`T=2\pi /\omega `$) is $`(2\sqrt{2})\pi `$ for each of the five geons.
### 3.5 Solution at the origin
The results of Section 3.4 suggest that in the neighborhood of the origin
$`R`$ $`=`$ $`R_0+R_{\mathrm{}}\mathrm{ln}r`$
$`u`$ $`=`$ $`u_0`$
$`v`$ $`=`$ $`v_2r^2,`$ (82)
where $`R_0`$, $`R_{\mathrm{}}`$, $`u_0`$ and $`v_2`$ are real constants and $`u_0,v_2>0`$. These expressions satisfy Eqs. (73) through (75) as $`r0`$ provided
$$R_{\mathrm{}}=\pm \sqrt{2}.$$
(83)
Now that we know the behavior of the fields as $`r0`$ we are in a position to justify two earlier assumptions. From expression (35) for $`\stackrel{~}{Q}^2(r)`$ we obtain
$$\underset{r0}{lim}\stackrel{~}{Q}^2(r)=\underset{r0}{lim}\left[\frac{r\mathrm{\Phi }^{}(r)}{\sqrt{u_0v_2}}\right]=0,$$
(84)
justifying the assumption $`\stackrel{~}{Q}^2(0)=0`$. And from expression (50) for $`\stackrel{~}{M}(r)`$ we obtain
$$\underset{r0}{lim}\stackrel{~}{M}(r)=\underset{r0}{lim}\left[\frac{ru^{}(r)}{2\sqrt{u_0v_2}}\right]=0,$$
(85)
justifying the assumption $`\stackrel{~}{M}(0)=0`$.
The wave function $`R`$ diverges logarithmically at the origin. This singularity is weak enough that terms proportional to $`R^2`$ in the components of the energy-momentum tensor (17) are integrable. But terms of the form $`(R^{})^2/v`$ diverge as $`r^4`$ and are not integrable. These nonintegrable terms cancel out of the expression $`2๐ณ_0{}_{}{}^{0}๐ณ`$ so the total energy is finite.
Substituting expressions (82) into Eq. (18) for $`๐ณ`$ and using $`๐ฑ=8\pi ๐ณ`$ gives
$$๐ฑ=\frac{2}{v_2r^4}$$
(86)
for the dominant behavior of the scalar curvature as $`r0`$. Thus the scalar curvature diverges at the center of each geon.
## 4 Implications of infinite curvature
Because the scalar curvature diverges at $`r=0`$ it is not possible to establish a locally flat coordinate system there. So the locus $`r=0`$ must be excluded from spacetime, leaving a hole in the spacetime manifold. The presence of this spacetime singularity raises the following questions:
1. Do physical parameters (such as mass and charge) diverge?
2. Do arbitrary boundary conditions arise at the singularity?
3. Can particles and photons encounter the singularity and, if so, what happens to them?
Questions 1 and 2 have already been addressed: the geon mass, charge, angular momentum and action are all finite, and the wave function is normalizable; and the conditions $`\stackrel{~}{Q}^2(0)=0`$ and $`\stackrel{~}{M}(0)=0`$ are not arbitrary boundary conditions at the singularityโthey are consequences of the model \[see Eqs. (84) and (85)\].
The answer to question 3 is not so clear-cut. In a classical (non-quantum) theory one considers spacetime pathological if the world line of a freely falling test particle (a point particle of negligible mass) does not exist after (or before) a finite interval of proper time. Such pathological spacetimes are said to be timelike geodesically incomplete In a quantum theory, however, no real particle can serve as a test particle in the neighborhood of a singularity, since such a particle would have to be so small (short wavelength, large energy) that it would itself dominate local spacetime structure. Perhaps the concept of geodesic completeness can be extended to the quantum realm , but at the moment no consensus exists on how to do this.
So we will not be able to answer question 3 definitively. But we will present a classical analysis which suggests the singularity is benign. In sections 4.1 and 4.2 we calculate timelike and null geodesicsโthe paths of freely falling particles and photonsโin the geon spacetime. We will find that the geon core is sufficiently repulsive to cloak the singularity and leave the spacetime timelike and null geodesically complete.
### 4.1 Geodesic equations
Consider moving along a geodesic with velocity $`๐^\mu `$. The trajectory obeys
$$\frac{d๐_\alpha }{ds}=\frac{1}{2}๐_{\mu \nu ,\alpha }๐^\mu ๐^\nu $$
(87)
subject to the constraint
$$๐^\mu ๐_\mu =\zeta ,$$
(88)
where $`\zeta =0,1`$ for null and timelike geodesics, and $`s`$ (proper time for a timelike geodesic) parameterizes the trajectory. $`๐_{\mu \nu }`$ is independent of $`t`$ and $`\varphi `$ so equation (87) implies $`๐_t`$ and $`๐_\varphi `$ are constant along the trajectory. Given the spherical symmetry we can (without loss of generality) consider trajectories confined to the equatorial plane ($`\theta =\pi /2`$, $`๐_\theta =๐^\theta =0`$). Thus
$$๐_\mu =\left(\begin{array}{c}\epsilon \\ v\dot{r}\\ 0\\ \mathrm{}\end{array}\right)\text{ and }๐^\mu =\left(\begin{array}{c}\epsilon /u\\ \dot{r}\\ 0\\ \mathrm{}/r^2\end{array}\right),$$
(89)
where $`\dot{r}dr/ds`$, and $`\epsilon `$ and $`\mathrm{}`$ are real constants. Then equation (88) gives
$$uv\dot{r}^2=\epsilon ^2V^2,$$
(90)
where
$$V^2u\left(\zeta +\frac{\mathrm{}^2}{r^2}\right).$$
(91)
The evolution of $`\varphi `$ is determined by
$$\dot{\varphi }d\varphi /ds=v^\varphi =\mathrm{}/r^2,$$
(92)
and this equation and (90) completely determine the trajectory. Since $`u,v0`$ equation (90) can be satisfied only when $`\epsilon ^2V^2`$, and the condition $`\epsilon ^2=V^2`$ corresponds to a turning point of the radial motion. We will refer to $`V^2`$ as the effective potential (for radial motion). The effective potentials for timelike and null geodesics in the spacetime surrounding the $`n=0`$ geon are plotted in Figs. 7 and 8.
### 4.2 Geodesics near the singularity
As $`r0`$ the metric approaches
$$ds^2=u_0dt^2v_2r^2dr^2r^2d\theta ^2r^2\mathrm{sin}^2\theta d\varphi ^2.$$
(93)
Using the results of section 4.1 we will explore the equatorial geodesics of this metric.
Let $`r_0`$ represent the minimum value of $`r`$ along a geodesic. By setting $`\epsilon ^2=V^2`$ we find
$$r_0=\sqrt{\frac{u_0\mathrm{}^2}{\epsilon ^2u_0\zeta }}.$$
(94)
The geodesic trajectory is determined by
$$\frac{dr}{d\varphi }=\frac{\dot{r}}{\dot{\varphi }}=\pm \frac{1}{\sqrt{v_2}}\sqrt{\left(\frac{r}{r_0}\right)^21}$$
(95)
which can be integrated to give
$$\varphi =\pm r_0^{}\mathrm{ln}\left[\frac{r^{}}{r_0^{}}+\sqrt{\left(\frac{r^{}}{r_0^{}}\right)^21}\right],$$
(96)
where we have introduced the scaled radial coordinate $`r^{}=\sqrt{v_2}r`$ and chosen the constant of integration so $`\varphi =0`$ when $`r^{}=r_0^{}`$. Geodesics corresponding to (96) are plotted in terms of the Cartesian coordinates
$`x`$ $``$ $`r^{}\mathrm{cos}\varphi `$
$`y`$ $``$ $`r^{}\mathrm{sin}\varphi `$ (97)
in Fig. 9. Test particles are attracted when $`r^{}1`$ and repelled when $`r^{}1`$. In the limit of zero angular momentum ($`\mathrm{}=r_0^{}=0`$) the geodesic is coincident with the positive $`x`$ axis, corresponding to a test particle rebounding directly backwards.
Now let us calculate the behavior of $`s`$ as a function of $`r`$ for particles and photons following the geodesics of Fig. 9. From Eq. (90) we obtain
$$\left(r\dot{r}\right)^2=\frac{\epsilon ^2V^2}{u_0v_2}=\kappa ^2\left[1\left(\frac{r_0}{r}\right)^2\right],$$
(98)
where
$$\kappa \sqrt{\frac{\epsilon ^2u_0\zeta }{u_0v_2}}.$$
(99)
Integrating (98) we obtain
$$s^{}=\pm \frac{(r_0^{})^2}{2}\left\{\frac{r^{}}{r_0^{}}\sqrt{\left(\frac{r^{}}{r_0^{}}\right)^21}+\mathrm{ln}\left[\frac{r^{}}{r_0^{}}+\sqrt{\left(\frac{r^{}}{r_0^{}}\right)^21}\right]\right\},$$
(100)
where we have introduced the scaled parameter $`s^{}=v_2\kappa s`$ and chosen the constant of integration so $`s^{}=0`$ when $`r^{}=r_0^{}`$. Equation (100) is plotted in Fig. 10 for various values of $`r_0^{}`$.
Figures 9 and 10 suggest that the core of the geon behaves like a repulsive potential barrier, effectively cloaking the singularity. Freely falling test particles with zero angular momentum ($`\mathrm{}=r_0^{}=0`$) do encounter the singularity, but they reboundโtheir world lines do not vanish. The spacetime is timelike and null geodesically complete.
One may object that any test particle encountering the singularity will be destroyed by infinite tidal forces. But keep in mind that no real particle can serve as a test particle at the Planck scale, so infinite tidal forces at the singularity do not necessarily pose a physical conundrum.
## 5 Discussion
The use of boundary condition 2 to force correspondence between Eqs. (61) and (63) can be motivated in the following way. Consider two different localized physical systems A and B with zero angular momentum and identical charges and masses. Suppose A is accurately described by the flat spacetime Klein-Gordon equation (59) while B is described by the curved spacetime equation (22). Viewed from a sufficiently large distance each system appears as a point particle with a wave function given by (60). Using local measurements of the probability density $`\overline{\mathrm{\Psi }}\mathrm{\Psi }`$ an asymptotic observer could determine $`\sigma `$ and thereby discriminate between the particles, despite the fact that both have identical charge, angular momentum and mass. It is physically appealing to expect the asymptotic behavior of $`\overline{\mathrm{\Psi }}\mathrm{\Psi }`$ to be uniquely determined by these three parameters. Indeed, this is what we typically mean by the word โparticleโ. By selecting $`\omega `$ to force $`\sigma `$ to be identical for particles A and B we guarantee that the asymptotic behavior of $`\overline{\mathrm{\Psi }}\mathrm{\Psi }`$ depends only on charge, angular momentum and mass, and is independent of internal details.
At the Planck scale many physicists expect classical notions of spacetime to failโa concept conveyed by the phrase (also coined by Wheeler) โspacetime foamโ. A successful quantum theory of gravity would, it is thought, flesh out the details of spacetime foam and erase the singularities associated with point particles.
The geon model developed here suggests a different perspective. Classical spacetime is assumed valid at the Planck scale and point particles are replaced by eignemodes of a quantum field. Singularities reminiscent of point particles remain, but they do not disturb the geodesic completeness of spacetime. Multi-particle systems would, presumably, correspond to multiple excitations of the geon field and all particle interactions (including those involving wave-packet reduction) would ultimately derive from the action (1). In short, the geon perspective replaces the search for a โquantum theory of gravityโ with the search for a โgravitational theory of quantaโ.
Some questions for future investigation come to mind:
1. Do charged solutions of Eq. (69) exist?
2. Do solutions with nonzero angular momentum exist? The spherically symmetric trial solution \[Eqs. (13), (14) and (15)\] could be replaced by one with axial symmetry. However, the trial solution would now involve $`\theta `$ and $`\varphi `$ as well as $`r`$ and $`t`$, and the metric tensor and vector potential would have additional nonzero components, so the solution would be much more challenging. Furthermore, it is not obvious what functional form should replace metric tensor (13).
3. What is the physical significance of the angular frequency $`\omega =M/\sqrt{2}`$?
4. What is the physical significance of the stationary value of the action and the action per cycle?
5. The action (1) is, arguably, the simplest which includes gravitation, electromagnetism and quantum mechanics. But there is no compelling reason to regard it as correct. \[8, p. 144\] It would be interesting to see how other terms in the action (such as a cosmological constant or conformal coupling \[9, p. 116\]) affect geon solutions.
The geons described here are far too massive to correspond to any known particle. They would interact gravitationally with ordinary matter so they appear to be candidates for dark matter. A number of workers have considered Planck-mass particles as dark matter (see and references therein) but it is not clear whether such models can be successfully incorporated into standard cosmology. The density of dark matter within a galactic halo is thought to be about $`0.3\text{GeV}\text{cm}^3`$ so, if all dark matter is composed of Plank-mass ($`1.2\times 10^{19}\text{GeV}`$) geons, the local geon number density is about $`0.3\times 10^{19}\text{cm}^3`$. This tiny density and weak coupling to ordinary matter would make the detection of such particles difficult.
## Appendix A Notation and conventions
Our notation follows Dirac . Greek indices take on the values $`0,1,2,3`$ and repeated indices are summed. The spacetime coordinates are $`๐^\mu `$ with $`๐^0=t`$ and the metric signature is $`+`$$``$$``$$``$.
The curvature tensor is
$$๐ฑ_{}^{\alpha }{}_{\mu \nu \beta }{}^{}\mathsf{\Gamma }_{}^{\alpha }{}_{\mu \nu ,\beta }{}^{}+\mathsf{\Gamma }_{}^{\alpha }{}_{\mu \beta ,\nu }{}^{}\mathsf{\Gamma }_{}^{\sigma }{}_{\mu \nu }{}^{}\mathsf{\Gamma }_{}^{\alpha }{}_{\sigma \beta }{}^{}+\mathsf{\Gamma }_{}^{\sigma }{}_{\mu \beta }{}^{}\mathsf{\Gamma }_{}^{\alpha }{}_{\sigma \nu }{}^{},$$
(101)
the Ricci tensor is $`๐ฑ_{\mu \nu }๐ฑ_{}^{\alpha }{}_{\mu \nu \alpha }{}^{}`$ and the scalar curvature is $`๐ฑ๐ฑ_{\mu }^{}{}_{}{}^{\mu }`$, where
$$\mathsf{\Gamma }_{\alpha \mu \nu }\frac{1}{2}\left(๐_{\alpha \mu ,\nu }+๐_{\alpha \nu ,\mu }๐_{\mu \nu ,\alpha }\right)$$
(102)
is the Christoffel symbol and a comma preceding some lower index $`\mu `$ denotes partial differentiation with respect to $`๐^\mu `$. The metric tensor $`๐_{\mu \nu }`$ is symmetric, its contraction is
$$๐_{\mu \nu }๐^{\mu \nu }=๐_{\mu }^{}{}_{}{}^{\mu }=4,$$
(103)
and, for notational brevity,
$$\sqrt{|det\left(๐_{\mu \nu }\right)|}.$$
(104)
A semicolon or colon preceding some lower index $`\mu `$ denotes a covariant derivative with respect to $`๐^\mu `$. A semicolon denotes the covariant derivative of general relativity. For a scalar, covariant vector and contravariant vector
$`\mathrm{\Psi }_{;\mu }`$ $``$ $`\mathrm{\Psi }_{,\mu }`$
$`๐ต_{\mu ;\nu }`$ $``$ $`๐ต_{\mu ,\nu }\mathsf{\Gamma }_{}^{\alpha }{}_{\mu \nu }{}^{}๐ต_\alpha `$
$`๐ต_{}^{\mu }{}_{;\nu }{}^{}`$ $``$ $`๐ต_{}^{\mu }{}_{,\nu }{}^{}+\mathsf{\Gamma }_{}^{\mu }{}_{\alpha \nu }{}^{}๐ต^\alpha `$ (105)
and the procedure generalizes in the usual way to tensors of any rank and mixture of covariant and contravariant indices. A colon denotes the generalized covariant derivative
$$๐ด_{:\mu }๐ด_{;\mu }+iQ๐ _\mu ๐ด$$
(106)
for an arbitrary tensor $`๐ด`$.
## Appendix B Asymptotic wave function
In the asymptotic regime \[see Eqs. (53) and (56)\]
$$\mathrm{\Phi }=Q^2r^1+O(r^2)$$
(107)
and
$$u=12Mr^1+Q^2r^2+O(r^3).$$
(108)
By expanding $`u^1`$ in terms of $`r^1`$ as $`r\mathrm{}`$ we obtain
$$v=1+2Mr^1+\left(4M^2Q^2\right)r^2+O(r^3).$$
(109)
Consider a trial solution for $`R`$ of the form
$$R=R_{\mathrm{}}\left[1+ar^1+O(r^2)\right]r^{1\sigma }e^{kr},$$
(110)
where $`a`$ is a real constant. Substituting the above expressions into Eq. (22) and expanding in terms of $`r^1`$ as $`r\mathrm{}`$ gives
$$0=\kappa +\left[\kappa a2M\left(M^22\omega ^2\right)2Q^2\omega +2k\sigma \right]r^1+O(r^2),$$
(111)
where
$$\kappa k^2+\omega ^2M^2.$$
(112)
Equation (111) must be satisfied term-by-term so $`\kappa =0`$ and we obtain Eqs. (62) and (63).
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